diff --git a/.gitattributes b/.gitattributes index 2366375bd4d80f4d0b7a823227242d5beb3e55cc..e08d1a5e054e43713e24fe2513d909cb7deea08f 100644 --- a/.gitattributes +++ b/.gitattributes @@ -216,3 +216,4 @@ data_all_eng_slimpj/shuffled/split/split_finalad/part-01.finalad filter=lfs diff data_all_eng_slimpj/shuffled/split/split_finalac/part-14.finalac filter=lfs diff=lfs merge=lfs -text data_all_eng_slimpj/shuffled/split/split_finalac/part-02.finalac filter=lfs diff=lfs merge=lfs -text data_all_eng_slimpj/shuffled/split/split_finalac/part-05.finalac filter=lfs diff=lfs merge=lfs -text +data_all_eng_slimpj/shuffled/split/split_finalac/part-10.finalac filter=lfs diff=lfs merge=lfs -text diff --git a/data_all_eng_slimpj/shuffled/split/split_finalac/part-10.finalac b/data_all_eng_slimpj/shuffled/split/split_finalac/part-10.finalac new file mode 100644 index 0000000000000000000000000000000000000000..3a0acd840070bcf39703f912e9f4fec009fc3af9 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split/split_finalac/part-10.finalac @@ -0,0 +1,3 @@ +version https://git-lfs.github.com/spec/v1 +oid sha256:61bb07281110301fbe807b058ccc9355379aad39710706a340006cd853fc9696 +size 12576663155 diff --git a/data_all_eng_slimpj/shuffled/split2/finalzbvb b/data_all_eng_slimpj/shuffled/split2/finalzbvb new file mode 100644 index 0000000000000000000000000000000000000000..7f18345666623ef194e261a2839acef8cb16e968 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzbvb @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\nLearning to learn is essential in human intelligence but is still a wide area of research in machine learning.\n\\textit{Meta-learning} has emerged as a popular approach to enable models to perform well on new tasks using limited data.\nIt involves first a \\textit{meta-training} process, when the model learns valuable features from a set of tasks.\nThen, at test time, using only few datapoints from a new, unseen task, the model (1) \\textit{adapts} to this new task (i.e., performs \\textit{few-shot learning} with \\textit{context data}), \nand then (2) \\textit{infers} by making predictions on new, unseen \\textit{query inputs} from the same task.\nA popular baseline for meta-learning, which has attracted a large amount of attention, is Model-Agnostic Meta-Learning (MAML) \\citep{maml}, in which the adaptation process consists of fine-tuning the parameters of the model via gradient descent.\n\nHowever, meta-learning methods can often struggle in several ways when deployed in challenging real-world scenarios. First, when context data is too limited to fully identify the test-time task, accurate prediction can be challenging. As these predictions can be untrustworthy, this necessitates the development of meta-learning methods that can express uncertainty during adaptation \\citep{bayesian_maml, alpaca}. In addition, meta-learning models may not successfully adapt to ``unusual'' tasks, i.e., when test-time context data is drawn from an \\textit{out-of-distribution} (OoD) task not well represented in the training dataset \\citep{ood_maml, meta_learning_ood}.\nFinally, special care has to be taken when learning tasks that have a large degree of heterogeneity.\nAn important example is the case of tasks with a \\textit{multimodal} distribution, i.e., when there are no common features shared across all the tasks, but the tasks can be broken down into subsets (modes) in a way that the ones from the same subset share common features \\citep{mmaml}.\n\n\\textbf{Our contributions.}~\\, We present \\textsc{UnLiMiTD}{} (\\textit{uncertainty-aware meta-learning for multimodal task distributions}), a novel meta-learning method that leverages probabilistic tools to address the aforementioned issues.\nSpecifically, \\textsc{UnLiMiTD}{} models the true distribution of tasks with a learnable distribution constructed over a linearized neural network and uses analytic Bayesian inference to perform uncertainty-aware adaption.\nWe present three variants (namely, \\approach-\\textsc{I}, \\approach-\\textsc{R}{}, and \\approach-\\textsc{F}) that reflect a trade-off between learning a rich prior distribution over the weights and maintaining the full expressivity of the network; we show that \\approach-\\textsc{F}{} strikes a balance between the two, making it the most appealing variant.\nFinally, we demonstrate that (1) our method allows for efficient probabilistic predictions on in-distribution tasks, that compare favorably to, and in most cases outperform, the existing baselines, (2) it is effective in detecting context data from OoD tasks at test time, and that (3) both these findings continue to hold in the multimodal task-distribution setting.\n\nThe rest of the paper is organized as follows. Section~\\ref{sec:problem_statement} formalizes the problem. Section~\\ref{sec:background} presents background information on the linearization of neural networks and Bayesian linear regression.\nWe detail our approach and its three variants in Section~\\ref{sec:approach}.\nWe discuss related work in detail in Section~\\ref{sec:related_work}. Finally, we present our experimental results concerning the performance of \\textsc{UnLiMiTD}{} in Section~\\ref{sec:results} and conclude in Section~\\ref{sec:conclusion}.\n\n\n\n\n\\begin{figure}[t]\n \\centering\n \\includegraphics[width=\\textwidth]{images\/flowchart.pdf}\n \\caption{ The true task distribution $p(f)$ can be multimodal, i.e., containing multiple clusters of tasks (e.g., lines and sines). Our approach \\textsc{UnLiMiTD}{} fits $p(f)$ with a parametric, tuneable distribution $\\Tilde{p}_\\xi(f)$ yielded by Bayesian linear regression on a linearized neural network.}\n \\label{fig:flowchart}\n\\end{figure}\n\n\\section{Problem statement}\\label{sec:problem_statement}\nA task $\\mathcal{T}^i$ consists of a function $f_i$ from which data is drawn.\nAt test time, the prediction steps are broken down into (1) \\textit{adaptation}, that is identifying $f_i$ using $K$ context datapoints $(\\xinput^i, \\youtput^i)$ from the task, and (2) \\textit{inference}, that is making predictions for $f_i$ on the \\textit{query inputs} $\\xinput^i_*$.\nLater the predictions can be compared with the \\textit{query ground-truths} $\\youtput^i_*$ to estimate the quality of the prediction, for example in terms of mean squared error (MSE).\nThe meta-training consists in learning valuable features from a \\textit{cluster of tasks}, which is a set of similar tasks (e.g., sines with different phases and amplitudes but same frequency), so that at test time the predictions can be accurate on tasks from the same cluster.\nWe take a probabilistic, functional perspective and represent a cluster by $p(f)$, a theoretical distribution over the function space that describes the probability of a task belonging to the cluster.\nLearning $p(f)$ is appealing, as it allows for performing OoD detection in addition to making predictions. Adaptation amounts to computing the conditional distribution given test context data, and one can obtain an uncertainty metric by evaluating the negative log-likelihood (NLL) of the context data under $p(f)$.\n\nThus, our goal is to construct a parametric, learnable functional distribution $\\Tilde{p}_\\xi(f)$ that approaches the theoretical distribution $p(f)$, with a structure that allows tractable conditioning and likelihood computation, even in deep learning contexts.\nIn practice, however, we are not given $p(f)$, but only a meta-training dataset $\\mathcal{D}$ that we assume is sampled from $p(f)$: $\\mathcal{D}=\\{ (\\widetilde{\\xinput}^i, \\widetilde{\\youtput}^i)\\}_{i=1}^{N}$, where $N$ is the number of tasks available during training, and $(\\widetilde{\\xinput}^i, \\widetilde{\\youtput}^i) \\sim \\mathcal{T}^i$ is the entire pool of data from which we can draw subsets of context data $(\\xinput^i, \\youtput^i)$.\nConsequently, in the meta-training phase, we aim to optimize $\\Tilde{p}_\\xi(f)$ to capture properties of $p(f)$, using only the samples in $\\mathcal{D}$.\n\nOnce we have $\\Tilde{p}_\\xi(f)$, we can evaluate it both in terms of how it performs for few-shot learning (by comparing the predictions with the ground truths in terms of MSE), as well as for OoD detection \n(by measuring how well the NLL of context data serves to classify in-distribution tasks against OoD tasks, measured via the AUC-ROC score). \\nopagebreak[4]\n\n\n\n\n\\section{Background}\n\\label{sec:background}\n\\subsection{Bayesian linear regression and Gaussian Processes}\n\\label{sec:reglin}\nEfficient Bayesian meta-learning requires a tractable inference process at test time. In general, this is only possible analytically in a few cases.\nOne of them is the Bayesian linear regression with Gaussian noise and a Gaussian prior on the weights. \nViewing it from a nonparametric, functional approach, this model is equivalent to a Gaussian process (GP) \\citep{rasmussen}.\n\nLet ${\\bm{X}} = ({\\bm{x}}_1, \\dots, {\\bm{x}}_K) \\in \\mathbb{R}^{{N_x} \\times K}$ be a batch of $K$ ${N_x}$-dimensional inputs, and let ${\\bm{y}} = ({\\bm{y}}_1, \\dots, {\\bm{y}}_K) \\in \\mathbb{R}^{{N_y} K}$ be a vectorized batch of ${N_y}$-dimensional outputs. In the Bayesian linear regression model, these quantities are related according to\n$\n {\\bm{y}} = \\phi({\\bm{X}})^\\top \\hat{\\param} + \\varepsilon \\in \\mathbb{R}^{{N_y} K}\n$\nwhere $\\hat{\\param} \\in \\mathbb{R}^P$ are the weights of the model, and the inputs are mapped via $\\phi:\\mathbb{R}^{{N_x} \\times K} \\rightarrow \\mathbb{R}^{P \\times {N_y} K}$. Notice how this is a generalization of the usual one-dimensional linear regression (${N_y}=1$).\n\n\nIf we assume a Gaussian prior on the weights $\\hat{\\param} \\sim \\mathcal{N}({\\bm{\\mu}}, {\\bm{\\Sigma}})$ and a Gaussian noise $\\varepsilon \\sim \\mathcal{N}(\\bm{0}, {\\bm{\\Sigma}}_\\varepsilon)$ with ${\\bm{\\Sigma}}_\\varepsilon = \\sigma_\\varepsilon^2 {\\bm{I}}$, then the model describes a multivariate Gaussian distribution on ${\\bm{y}}$ for any ${\\bm{X}}$. Equivalently, this means that this model describes a GP distribution over functions, with mean and covariance function (or kernel)\n\\begin{align}\n\\begin{split}\n\\label{eq:prior_pred_dist}\n{\\bm{\\mu}}_{\\text{prior}} ({\\bm{x}}_t) & = \\phi({\\bm{x}}_t)^\\top {\\bm{\\mu}}, \\\\\n\\text{cov}_{\\text{prior}} ({\\bm{x}}_{t_1}, {\\bm{x}}_{t_2}) & = \\phi({\\bm{x}}_{t_1})^\\top {\\bm{\\Sigma}} \\phi({\\bm{x}}_{t_2}) + {\\bm{\\Sigma}}_\\varepsilon =: k_{\\bm{\\Sigma}}({\\bm{x}}_{t_1}, {\\bm{x}}_{t_2}) + {\\bm{\\Sigma}}_\\varepsilon .\n\\end{split}\n\\end{align}\nThis GP enables tractable computation of the likelihood of any batch of data $({\\bm{X}}, {\\bm{Y}})$ given this distribution over functions. The structure of this distribution is governed by the feature map $\\phi$ and the prior over the weights, specified by ${\\bm{\\mu}}$ and ${\\bm{\\Sigma}}$.\n\nThis distribution can also easily be conditioned to perform inference. Given a batch of data $({\\bm{X}}, {\\bm{Y}})$, the posterior predictive distribution is also a GP, with an updated mean and covariance function\n\\begin{align}\n \\label{eq:post_pred_dist}\n \\begin{split}\n {\\bm{\\mu}}_{\\text{post}} ({\\bm{x}}_{t_*}) & = k_{\\bm{\\Sigma}}({\\bm{x}}_{t_*}, {\\bm{X}}) \\left( k_{\\bm{\\Sigma}}({\\bm{X}}, {\\bm{X}}) + {\\bm{\\Sigma}}_\\varepsilon \\right)^{-1} {\\bm{Y}}, \\\\\n \\text{cov}_{\\text{post}} ({\\bm{x}}_{{t_1}_*}, {\\bm{x}}_{{t_2}_*}) & = k_{\\bm{\\Sigma}}({\\bm{x}}_{{t_1}_*}, {\\bm{x}}_{{t_2}_*}) - k_{\\bm{\\Sigma}}({\\bm{x}}_{{t_1}_*}, {\\bm{X}}) \\left( k_{\\bm{\\Sigma}}({\\bm{X}}, {\\bm{X}}) + {\\bm{\\Sigma}}_\\varepsilon \\right)^{-1} k_{\\bm{\\Sigma}}({\\bm{X}}, {\\bm{x}}_{{t_2}_*}).\n \\end{split}\n\\end{align}\nHere, ${\\bm{\\mu}}_{\\text{post}}({\\bm{X}}_*)$ represents our model's adapted predictions for the test data, which we can compare to ${\\bm{Y}}_*$ to evaluate the quality of our predictions, for example, via mean squared error (assuming that test data is clean, following \\citet{rasmussen}). \nThe diagonal of $\\text{cov}_{\\text{post}}({\\bm{X}}_*, {\\bm{X}}_*)$ can be interpreted as a per-input level of confidence that captures the ambiguity in making predictions with only a limited amount of context data.\n\n\\subsection{The linearization of a neural network yields an expressive linear regression model}\n\\label{sec:linearization}\nAs discussed, the choice of feature map $\\phi$ plays an important role in specifying a linear regression model.\nIn the deep learning context, recent work has demonstrated that the linear model obtained when linearizing a deep neural network with respect to its weights at initialization, wherein the Jacobian of the network operates as the feature map, can well approximate the training behavior of wide nonlinear deep neural networks \\citep{jacot,nonlinear,liu2020linearity,shallow_nns_infinite_width,dnns_infinite_wdith}.\n\n\nLet $f$ be a neural network $f: \\left({\\bm{\\theta}}, {\\bm{x}}_t \\right) \\mapsto {\\bm{y}}_t$, where ${\\bm{\\theta}} \\in \\mathbb{R}^{P}$ are the parameters of the model, ${\\bm{x}} \\in \\mathbb{R}^{{N_x}}$ is an input and ${\\bm{y}} \\in \\mathbb{R}^{{N_y}}$ an output.\nThe linearized network (w.r.t. the parameters) around $\\param_0$ is\n\\begin{displaymath}\n f({\\bm{\\theta}}, {\\bm{x}}_t) - f(\\param_0, {\\bm{x}}_t) \\approx {\\bm{J}}_{\\bm{\\theta}}( f )(\\param_0, {\\bm{x}}_t) ({\\bm{\\theta}} - \\param_0),\n\\end{displaymath}\nwhere ${\\bm{J}}_{\\bm{\\theta}}(f)(\\cdot, \\cdot): \\mathbb{R}^P \\times \\mathbb{R}^{N_x} \\rightarrow \\mathbb{R}^{{N_y} \\times P}$ is the Jacobian of the network (w.r.t. the parameters).\n\nIn the case where the model accepts a batch of $K$ inputs ${\\bm{X}} = ({\\bm{x}}_1, \\dots, {\\bm{x}}_K)$ and returns ${\\bm{Y}} = ({\\bm{y}}_1, \\dots, {\\bm{y}}_K)$, we generalize $f$ to $g: \\mathbb{R}^P \\times \\mathbb{R}^{{N_x} \\times K} \\rightarrow \\mathbb{R}^{{N_y} \\times K}$, with ${\\bm{Y}} = g({\\bm{\\theta}}, {\\bm{X}})$.\nConsequently, we generalize the linearization:\n\\begin{displaymath}\ng({\\bm{\\theta}}, {\\bm{X}}) - g(\\param_0, {\\bm{X}}) \\approx {\\bm{J}}(\\param_0, {\\bm{X}}) ({\\bm{\\theta}} - \\param_0),\n\\end{displaymath}\nwhere ${\\bm{J}}(\\cdot, \\cdot): \\mathbb{R}^P \\times \\mathbb{R}^{{N_x} \\times K} \\rightarrow \\mathbb{R}^{{N_y} K \\times P}$ is a shorthand for ${\\bm{J}}_{\\bm{\\theta}}(g)(\\cdot, \\cdot)$.\nNote that we have implicitly vectorized the outputs, and throughout the work, we will interchange the matrices $\\mathbb{R}^{{N_y} \\times K}$ and the vectorized matrices $\\mathbb{R}^{{N_y} K}$.\n\nThis linearization can be viewed as the ${N_y} K$-dimensional linear regression\n\\begin{equation}\n \\label{eq:linearized_network}\n {\\bm{z}} = \\phi_{\\param_0}({\\bm{X}})^\\top \\hat{\\param} \\in \\mathbb{R}^{{N_y} K},\n\\end{equation}\nwhere the feature map $\\phi_{\\param_0}(\\cdot): \\mathbb{R}^{{N_x} \\times K} \\rightarrow \\mathbb{R}^{P \\times {N_y} K}$ is the transposed Jacobian ${\\bm{J}}(\\param_0, \\cdot)^\\top$.\nThe parameters of this linear regression $\\hat{\\param} = \\left( {\\bm{\\theta}} - \\param_0 \\right)$ are the \\textit{correction} to the parameters chosen as the linearization point.\nEquivalently, this can be seen as a kernel regression with the kernel $ k_{\\param_0}({\\bm{X}}_1,{\\bm{X}}_2) = {\\bm{J}}(\\param_0, {\\bm{X}}_1) {\\bm{J}}(\\param_0, {\\bm{X}}_2)^\\top$, which is commonly referred to as the Neural Tangent Kernel (NTK) of the network. Note that the NTK depends on the linearization point $\\param_0$. \nBuilding on these ideas, \\citet{maddox} show that the NTK obtained via linearizing a DNN \\textit{after} it has been trained on a task yields a GP that is well-suited for adaptation and fine-tuning to new, similar tasks. Furthermore, they show that networks trained on similar tasks tend to have similar Jacobians, suggesting that neural network linearization can yield an effective model for multi-task contexts such as meta-learning. In this work, we leverage these insights to construct our parametric functional distribution $\\Tilde{p}_\\xi(f)$ via linearizing a neural network model.\n\n\\section{Our approach: \\textsc{UnLiMiTD}}\n\\label{sec:approach}\nIn this section, we describe our meta-learning algorithm \\textsc{UnLiMiTD}{} and the construction of a parametric functional distribution $\\Tilde{p}_\\xi(f)$ that can model the true underlying distribution over tasks $p(f)$.\nFirst, we focus on the single cluster case, where a Gaussian process structure on $\\Tilde{p}_\\xi(f)$ can effectively model the true distribution of tasks, and detail how we can leverage meta-training data $\\mathcal{D}$ from a single cluster of tasks to train the parameters $\\xi$ of our model.\nNext, we will generalize our approach to the multimodal setting, with more than one cluster of tasks. Here, we construct $\\Tilde{p}_\\xi(f)$ as a mixture of GPs and develop a training approach that can automatically identify the clusters present in the training dataset without requiring the meta-training dataset to contain any additional structure such as cluster labels.\n\n\\subsection{Tractably structuring the prior predictive distribution over functions via a Gaussian distribution over the weights}\nIn our approach, we choose $\\Tilde{p}_\\xi(f)$ to be the \nGP distribution over functions that arises from a Gaussian prior on the weights of the linearization of a neural network (\\eqref{eq:linearized_network}). Consider a particular task $\\mathcal{T}^i$ and a batch of $K$ context data $(\\xinput^i, \\youtput^i)$.\nThe resulting prior predictive distribution, derived from \\eqref{eq:prior_pred_dist} after evaluating on the context inputs, is ${\\bm{Y}} | \\xinput^i \\sim \\mathcal{N}( {\\bm{\\mu}}_{\\youtput \\mid \\xcontextinput}, {\\bm{\\Sigma}}_{\\youtput \\mid \\xcontextinput})$, where\n\\begin{equation}\n \\label{eq:prior_pred_dist_ntk}\n {\\bm{\\mu}}_{\\youtput \\mid \\xcontextinput} = {\\bm{J}}(\\param_0, \\xinput^i) {\\bm{\\mu}}, \\quad {\\bm{\\Sigma}}_{\\youtput \\mid \\xcontextinput} = {\\bm{J}}(\\param_0, \\xinput^i) {\\bm{\\Sigma}} {\\bm{J}}(\\param_0, \\xinput^i)^\\top + {\\bm{\\Sigma}}_\\varepsilon.\n\\end{equation}\nIn this setup, the parameters $\\xi$ of $\\Tilde{p}_\\xi(f)$ that we wish to optimize are the linearization point $\\param_0$, and the parameters of the prior over the weights $({\\bm{\\mu}}, {\\bm{\\Sigma}})$.\nGiven this Gaussian prior, it is straightforward to compute the joint NLL of the context labels $\\youtput^i$,\n\\begin{align}\n\\label{eq:single-nll}\n \\mathrm{NLL}(\\xinput^i, \\youtput^i) = \\frac12\\left( \\left\\| \\youtput^i - {\\bm{\\mu}}_{\\youtput \\mid \\xcontextinput} \\right\\|^2_{{\\bm{\\Sigma}}_{\\youtput \\mid \\xcontextinput}^{-1}} + \\log\\det {\\bm{\\Sigma}}_{\\youtput \\mid \\xcontextinput} + {N_y} K \\log 2 \\pi \\right).\n\\end{align}\nThe NLL (a) serves as a loss function quantifying the quality of $\\xi$ during training and (b) serves as an uncertainty signal at test time to evaluate whether context data $(\\xinput^i, \\youtput^i)$ is OoD.\nGiven this model, \\textit{adaptation} is tractable as we can condition this GP on the context data analytically. In addition, we can efficiently make probabilistic predictions by evaluating the mean and covariance of the resulting posterior predictive distribution on the query inputs, using \\eqref{eq:post_pred_dist}.\n\n\\subsubsection{Parameterizing the prior covariance over the weights}\n\\label{sec:prior_covariance}\nWhen working with deep neural networks, the number of weights $P$ can surpass $10^6$. While it remains tractable to deal with $\\param_0$ and ${\\bm{\\mu}}$, whose memory footprint grows linearly with $P$, it can quickly become intractable to make computations with (let alone store) a dense prior covariance matrix over the weights ${\\bm{\\Sigma}} \\in \\mathbb{R}^{P \\times P}$. Thus, we must impose some structural assumptions on the prior covariance to scale to deep neural network models.\n\n\n\\textbf{Imposing a unit covariance.}~\\, One simple way to tackle this issue would be to remove ${\\bm{\\Sigma}}$ from the learnable parameters $\\xi$, i.e., fixing it to be the identity ${\\bm{\\Sigma}} = {\\bm{I}}_{P}$. In this case, $\\xi = (\\param_0, {\\bm{\\mu}})$. \nThis computational benefit comes at the cost of model expressivity, as we lose a degree of freedom in how we can optimize our learned prior distribution $\\Tilde{p}_\\xi(f)$. In particular, we are unable to choose a prior over the weights of our model that captures correlations between elements of the feature map.\n\n\\textbf{Learning a low-dimensional representation of the covariance.}~\\,\nAn alternative is to learn a low-rank representation of ${\\bm{\\Sigma}}$, allowing for a learnable weight-space prior covariance that can encode correlations. Specifically, we consider a covariance of the form ${\\bm{\\Sigma}} = {\\bm{Q}}^\\top \\diag{{\\bm{s}}^2} {\\bm{Q}}$, where ${\\bm{Q}}$ is a fixed projection matrix on an $s$-dimensional subspace of $\\mathbb{R}^{P}$, while ${\\bm{s}}^2$ is learnable.\nIn this case, the parameters that are learned are $\\xi = (\\param_0, {\\bm{\\mu}}, {\\bm{s}})$.\nWe define ${\\bm{S}} := \\diag{{\\bm{s}}^2}$.\nThe computation of the covariance of the prior predictive (\\eqref{eq:prior_pred_dist_ntk}) could then be broken down into two steps:\n\\begin{displaymath}\n\\left\\{\n \\begin{array}{l}\n A := {\\bm{J}}(\\param_0, \\xinput^i) {\\bm{Q}}^\\top \\\\\n {\\bm{J}}(\\param_0, \\xinput^i) {\\bm{\\Sigma}} {\\bm{J}}(\\param_0, \\xinput^i)^\\top = A {\\bm{S}} A^\\top \n \\end{array}\n\\right.\n\\end{displaymath}\nwhich requires a memory footprint of $O(P(s + {N_y} K) )$, if we include the storage of the Jacobian.\nBecause ${N_y} K \\ll P$ in typical deep learning contexts, it suffices that $s \\ll P$ so that it becomes tractable to deal with this new representation of the covariance.\n\n\\textbf{A trade-off between feature-map expressiveness and learning a rich prior over the weights.} Note that even if a low-dimensional representation of ${\\bm{\\Sigma}}$ enriches the prior distribution over the weights, it also restrains the expressiveness of the feature map in the kernel by projecting the $P$-dimensional features ${\\bm{J}}(\\param_0, {\\bm{X}})$ on a subspace of size $s \\ll P$ via ${\\bm{Q}}$.\nThis presents a trade-off: we can use the full feature map, but limit the weight-space prior covariance to be the identity matrix by keeping ${\\bm{\\Sigma}} = {\\bm{I}}$ (case \\approach-\\textsc{I}). Alternatively, we could learn a low-rank representation of ${\\bm{\\Sigma}}$ by randomly choosing $s$ orthogonal directions in $\\mathbb{R}^{P}$, with the risk that they could limit the expressiveness of the feature map if the directions are not relevant to the problem that is considered (case \\approach-\\textsc{R}).\nAs a compromise between these two cases, we can choose the projection matrix more intelligently and project to the most impactful subspace of the full feature map --- in this way, we can reap the benefits of a tuneable prior covariance while minimizing the useful features that the projection drops. To select this subspace, we construct this projection map by choosing the top $s$ eigenvectors of the Fisher information matrix (FIM) evaluated on the training dataset $\\mathcal{D}$ (case \\approach-\\textsc{F}). Recent work has shown that the FIM for deep neural networks tends to have rapid spectral decay \\citep{scod}, which suggests that keeping only a few of the top eigenvectors of the FIM is enough to encode an expressive task-tailored prior. See Appendix~\\ref{app:fim} for more details. \n\n\\subsubsection{Generalizing the structure to a mixture of Gaussians}\n\\label{sec:mixture}\n\nWhen learning on multiple clusters of tasks, $p(f)$ can become non-unimodal, and thus cannot be accurately described by a single GP.\nInstead, we can capture this multimodality by structuring $\\Tilde{p}_\\xi(f)$ as a \\textit{mixture} of Gaussian processes.\n\n\\textbf{Building a more general structure.}~\\, We assume that at train time, a task $\\mathcal{T}^i$ comes from any cluster $\\left\\{\\mathcal{C}_j \\right\\}_{j=1}^{j=\\alpha}$ with equal probability.\nThus, we choose to construct $\\Tilde{p}_\\xi(f)$ as an equal-weighted mixture of $\\alpha$ Gaussian processes.\n\nFor each element of the mixture, the structure is similar to the single cluster case, where the parameters of the cluster's weight-space prior are given by $({\\bm{\\mu}}_j, {\\bm{\\Sigma}}_j)$. We choose to have both the projection matrix ${\\bm{Q}}$ and the linearization point $\\param_0$ (and hence, the feature map $\\phi(\\cdot) = {\\bm{J}}(\\param_0,\\cdot)$) shared across the clusters. This yields improved computational efficiency, as we can compute the projected features once, simultaneously, for all clusters.\nThis yields the parameters $\\xi_\\alpha = (\\param_0, {\\bm{Q}}, ({\\bm{\\mu}}_1, {\\bm{s}}_1), \\ldots, ({\\bm{\\mu}}_\\alpha, {\\bm{s}}_\\alpha))$.\n\nThis can be viewed as a mixture of linear regression models, with a common feature map but \nseparate, independent prior distributions over the weights for each cluster. These separate distributions are encoded using the low-dimensional representations ${\\bm{S}}_j$ for each ${\\bm{\\Sigma}}_j$. \nNotice how this is a generalization of the single cluster case, for when $\\alpha=1$, $\\Tilde{p}_\\xi(f)$ becomes a Gaussian and $\\xi_\\alpha = \\xi$\\footnote{In theory, it is possible to drop ${\\bm{Q}}$ and extend the identity covariance case to the multi-cluster setting; however, this leads to each cluster having an identical covariance function, and thus is not effective at modeling heterogeneous behaviors among clusters.}.\n\n\n\\textbf{Prediction and likelihood computation.}~\\, The NLL of a batch of inputs under this mixture model can be computed as\n\\begin{equation}\n \\label{eq:nll_mixture}\n \\mathrm{NLL}_{\\text{mixt}}(\\xinput^i, \\youtput^i) = \\log \\alpha - \\log \\add \\exp (-\\mathrm{NLL}_1(\\xinput^i, \\youtput^i), \\ldots, -\\mathrm{NLL}_\\alpha(\\xinput^i, \\youtput^i)),\n\\end{equation}\nwhere $\\mathrm{NLL}_j(\\xinput^i, \\youtput^i)$ is the NLL with respect to each individual Gaussian, as computed in \\eqref{eq:single-nll}, and $\\log\\add\\exp$ computes the logarithm of the sum of the exponential of the arguments, taking care to avoid underflow issues.\n\nTo make exact predictions, we would require conditioning this mixture model. As this is not directly tractable, we propose to first \\textit{infer the cluster} from which a task comes from, by identifying the Gaussian $\\mathcal{G}_{j_0}$ that yields the highest likelihood for the context data $\\left( \\xinput^i, \\youtput^i \\right)$. Then, we can \\textit{adapt} by conditioning $\\mathcal{G}_{j_0}$ with the context data and finally \\textit{infer} by evaluating the resulting posterior distribution on the queried inputs $\\xinput^i_*$.\n\n\\subsection{Meta-training the Parametric Task Distribution}\nThe key to our meta-learning approach is to estimate the quality of $\\Tilde{p}_\\xi(f)$ via the NLL of context data from training tasks, and use its gradients to update the parameters of the distribution $\\xi$.\nOptimizing this loss over tasks in the dataset draws $\\Tilde{p}_\\xi(f)$ closer to the empirical distribution present in the dataset, and hence towards the true distribution $p(f)$.\n\n\nWe present three versions of \\textsc{UnLiMiTD}, depending on the choice of structure of the prior covariance over the weights (see Section~\\ref{sec:prior_covariance} for more details).\n\\approach-\\textsc{I}{} (Algorithm~\\ref{alg:meta_training_identity}) is the meta-training with the fixed identity prior covariance. \\approach-\\textsc{R}{} and \\approach-\\textsc{F}{} (Algorithm~\\ref{alg:meta_training_learnt_cov}) learn a low-dimensional representation of that prior covariance, either with random projections or with FIM-based projections.\n\n\\begin{algorithm}[t]\n\\caption{\\footnotesize \\approach-\\textsc{I}: meta-training with identity prior covariance}\n\\footnotesize\n\\label{alg:meta_training_identity}\n\\begin{algorithmic}[1]\n \\State Initialize $\\param_0$, ${\\bm{\\mu}}$.\n \\ForAll{epoch}\n \\State Sample $n$ tasks $\\{ \\mathcal{T}^i, (\\xinput^i, \\youtput^i) \\}_{i=1}^{i=n}$\n \\ForAll{$\\mathcal{T}^i, (\\xinput^i, \\youtput^i)$}\n \\State $NLL_i \\gets \\Call{GaussNLL}{\\youtput^i; {\\bm{J}}{\\bm{\\mu}},~ {\\bm{J}}\\jac^\\top + {\\bm{\\Sigma}}_\\varepsilon}$ \\Comment{${\\bm{J}} = {\\bm{J}}(\\param_0, \\xinput^i)$}\n \\EndFor\n \\State Update $\\param_0$, ${\\bm{\\mu}}$ with $\\nabla_{\\param_0 \\cup {\\bm{\\mu}}} \\sum_i NLL_i$\n \\EndFor\n\\end{algorithmic}\n\\end{algorithm}\n\n\\begin{algorithm}[t]\n\\caption{\\footnotesize \\approach-\\textsc{R}{} and \\approach-\\textsc{F}{}: meta-training with a learnt covariance}\n\\footnotesize\n\\label{alg:meta_training_learnt_cov}\n\\begin{algorithmic}[1]\n \\If{using random projections}\n \\State Find random projection ${\\bm{Q}}$\n \\State Initialize $\\param_0$, ${\\bm{\\mu}}$, ${\\bm{s}}$\n \\ElsIf{using FIM-based projections}\n \\State Find intermediate $\\param_0$, ${\\bm{\\mu}}$ with \\approach-\\textsc{I}{} \\Comment{see Alg.~\\ref{alg:meta_training_identity}}\n \\State Find ${\\bm{Q}}$ via \\Call{FIMProj}{s}; initialize ${\\bm{s}}$. \\Comment{see Alg.~\\ref{alg:fim_proj}}\n \\EndIf\n \\ForAll{epoch}\n \\State Sample $n$ tasks $\\{ \\mathcal{T}^i, (\\xinput^i, \\youtput^i) \\}_{i=1}^{i=n}$\n \\ForAll{$\\mathcal{T}^i, (\\xinput^i, \\youtput^i)$}\n \\State $NLL_i \\gets \\Call{GaussNLL}{\\youtput^i; {\\bm{J}}{\\bm{\\mu}},~ {\\bm{J}} {\\bm{Q}}^\\top \\diag{{\\bm{s}}^2} {\\bm{Q}} {\\bm{J}}^\\top + {\\bm{\\Sigma}}_\\varepsilon}$ \\Comment{${\\bm{J}} = {\\bm{J}}(\\param_0, \\xinput^i)$}\n \\EndFor\n \\State Update $\\param_0$, ${\\bm{\\mu}}$, ${\\bm{s}}$ with $\\nabla_{\\param_0 \\cup {\\bm{\\mu}} \\cup {\\bm{s}}} \\sum_i NLL_i$\n \\EndFor\n\\end{algorithmic}\n\\end{algorithm}\n\n\\textbf{Computing the likelihood.}~\\,\nIn the algorithms, the function \\Call{GaussNLL}{$\\youtput^i$; $m$, $K$} stands for NLL of $\\youtput^i$ under the Gaussian $\\mathcal{N}(m, K)$ (see \\eqref{eq:single-nll}).\nIn the mixture case, we instead use \\Call{MixtNLL}{}, which wraps \\eqref{eq:nll_mixture} and calls \\Call{GaussNLL}{} for the individual NLL computations (see discussion in Section~\\ref{sec:mixture}).\nIn this case, ${\\bm{\\mu}}$ becomes $\\{{\\bm{\\mu}}_j\\}_{j=1}^{j=\\alpha}$ and ${\\bm{s}}$ becomes $\\{{\\bm{s}}_j\\}_{j=1}^{j=\\alpha}$ when applicable.\n\n\\textbf{Finding the FIM-based projections.}~\\,\nThe FIM-based projection matrix aims to identify the elements of $\\phi = {\\bm{J}}(\\param_0, {\\bm{X}})$ that are most relevant for the problem (see Section~\\ref{sec:prior_covariance} and Appendix~\\ref{app:fim}).\nHowever, this feature map evolves during training, because it is $\\param_0$-dependent.\nHow do we ensure that the directions we choose for ${\\bm{Q}}$ remain relevant during training?\nWe leverage results from \\citet{ntk_evolution}, stating that the NTK (the kernel associated with the Jacobian feature map, see Section~\\ref{sec:linearization}) changes significantly at the beginning of training and that its evolution slows down as training goes on. This suggests that as a heuristic, we can compute the FIM-based directions after partial training, as they are unlikely to deviate much after the initial training. \nFor this reason, \\approach-\\textsc{F}{}\n(Algorithm~\\ref{alg:meta_training_learnt_cov}) first calls \\approach-\\textsc{I}{} (Algorithm~\\ref{alg:meta_training_identity}) before computing the FIM-based ${\\bm{Q}}$ that yields intermediate parameters $\\param_0$ and ${\\bm{\\mu}}$.\nThen the usual training takes place with the learning of ${\\bm{s}}$ in addition to $\\param_0$ and ${\\bm{\\mu}}$.\n\n\\section{Related work}\n\\label{sec:related_work}\n\\textbf{Bayesian inference with linearized DNNs.}~\\,\nBayesian inference with neural networks is often intractable because the posterior predictive has rarely a closed-form expression.\nWhereas \\textsc{UnLiMiTD}{} linearizes the network to allow for practical Bayesian inference, existing work has used other approximations to tractably express the posterior.\nFor example, it has been shown that in the infinite-width approximation, the posterior predictive of a Bayesian neural network behaves like a GP \\citep{shallow_nns_infinite_width, dnns_infinite_wdith}. This analysis can in some cases yield a good approximation to the Bayesian posterior of a DNN \\citep{cnns_infinite_width}.\nIt is also common to use Laplace's method to approximate the posterior predictive by a Gaussian distribution and allow practical use of the Bayesian framework for neural networks.\nThis approximation relies in particular on the computation of the Hessian of the network: this is in general intractable, and most approaches use the so-called Gauss-Newton approximation of the Hessian instead \\citep{laplace_scalable}.\nRecently, it has been shown that the Laplace method using the Gauss-Newton approximation is equivalent to working with a certain linearized version of the network and its resulting posterior GP \\citep{laplace_linearization}.\n\nBayesian inference is applied in a wide range of subjects. For example, recent advances in transfer learning have been possible thanks to Bayesian inference with linearized neural networks.\n\\citet{maddox} have linearized pre-trained networks and performed domain adaptation by conditioning the prior predictive with data from the new task: the posterior predictive is then used to make predictions. Our approach leverages a similar adaption method and demonstrates how the prior distribution can be learned in a meta-learning setup.\n\n\\textbf{Meta-learning.}~\\,\nMAML is a meta-learning algorithm that uses as adaptation a few steps of gradient descent \\citep{maml}.\nIt has the benefit of being model-agnostic (it can be used on any model for which we can compute gradients w.r.t. the weights), whereas \\textsc{UnLiMiTD}{} requires the model to be a differentiable regressor.\nMAML has been further generalized to probabilistic meta-learning models such as PLATIPUS or BaMAML \\citep{bayesian_maml, probabilistic_maml}, where the simple gradient descent step is augmented to perform approximate Bayesian inference. These approaches, like ours, learn (during meta-training) and make use of (at test-time) a prior distribution on the weights. In contrast, however, \\textsc{UnLiMiTD}{} uses exact Bayesian inference at test-time.\nMAML has also been improved for multimodal meta-learning via MMAML \\citep{mmaml, revisit_mmaml}. Similarly to our method, they add a step to identify the cluster from which the task comes from \\citep{mmaml}.\nOoD detection in meta-learning has been studied by \\citet{ood_maml}, who build upon MAML to perform OoD detection in the classification setting, to identify unseen classes during training.\n\\citet{meta_learning_ood} also implemented OoD detection for classification, by learning a Gaussian mixture model on a latent space. \n\\textsc{UnLiMiTD}{} extends these ideas to the regression task, aiming to identify when test data is drawn from an unfamiliar function.\n\nALPaCA is a Bayesian meta-learning algorithm for neural networks, where only the last layer is Bayesian \\citep{alpaca}.\nSuch framework yields an exact linear regression that uses as feature map the activations right before the last layer.\nOur work is a generalization of ALPaCA, in the sense that \\textsc{UnLiMiTD}{} restricted to the last layer matches ALPaCA's approach.\nMore on this link between the methods is discussed in Appendix~\\ref{app:link_with_alpaca}.\n\n\\section{Results and discussion}\n\\label{sec:results}\nWe wish to evaluate four key aspects of \\textsc{UnLiMiTD}.\n(1) At test time, how do the probabilistic predictions compare to baselines?\n(2) How well does the detection of context data from OoD tasks perform?\n(3) How do these results hold in the multimodal setting?\n(4) Which approach performs better between (a) the identity covariance (\\approach-\\textsc{I}), (b) the low-dimensional covariance with random directions (\\approach-\\textsc{R}) and the compromise (c) using FIM-based directions (\\approach-\\textsc{F}) (see trade-off in Section~\\ref{sec:prior_covariance})?\nThat is, what is best between learning a rich prior distribution over the weights, keeping a full feature map, and a compromise between the two?\n\nWe consider a cluster of sine tasks, one of linear tasks and one of quadratic tasks, regression problems inspired from \\citet{mmaml}.\nDetails on the problems can be found in Appendix~\\ref{app:problem-details}.\n\n\\textbf{Unimodal meta-learning: The meta-learned prior accurately fits the tasks.}~\\,\nFirst, we investigate the performance of \\textsc{UnLiMiTD}{} on a unimodal task distribution consisting of sinusoids of varying amplitude and phase, using the single GP structure for $\\Tilde{p}_\\xi(f)$.\nWe compare the performance between \\approach-\\textsc{I}, \\approach-\\textsc{R}{} and \\approach-\\textsc{F}.\nWe also compare the results between training with an infinite amount of available sine tasks (infinite task dataset), and with a finite amount of available tasks (finite task dataset).\nMore training details can be found in Appendix~\\ref{app:train-details-single}.\nExamples of predictions at the test time are available in Figure~\\ref{fig:single-predictions}, along with confidence levels.\n\n\\begin{figure}[t]\n \\centering\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{images\/predictions\/single_pred_fim_infinite_1.pdf}\n \\end{subfigure}\n \\hfill\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{images\/predictions\/single_pred_fim_infinite_5.pdf}\n \\end{subfigure}\n \\hfill\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{images\/predictions\/single_pred_fim_infinite_10.pdf}\n \\end{subfigure}\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{images\/predictions\/single_pred_fim_finite_1.pdf}\n \\end{subfigure}\n \\hfill\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{images\/predictions\/single_pred_fim_finite_5.pdf}\n \\end{subfigure}\n \\hfill\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{images\/predictions\/single_pred_fim_finite_10.pdf}\n \\end{subfigure}\n \\caption{ Example of predictions for a varying number of context inputs $K$, after meta-training with \\approach-\\textsc{F}. Top: \\approach-\\textsc{F}, infinite task dataset. Bottom: \\approach-\\textsc{F}, finite task dataset. The standard deviation is from the posterior predictive distribution.\n Note how the uncertainty levels are coherent with the actual prediction error.\n Also, note how uncertainty decreases when there is more context data.\n Notice how \\approach-\\textsc{F}{} recovers the shape of the sine even with a low number of context inputs.\n Finally, note how \\approach-\\textsc{F}{} is able to reconstruct the sine even when trained on fewer tasks (bottom). More comprehensive plots available in Figure~\\ref{fig:single-predictions-full}.}\n \\label{fig:single-predictions}\n\\end{figure}\n\nIn both OoD detection and quality of predictions, \\approach-\\textsc{R}{} and \\approach-\\textsc{F}{} \nperform better than \\approach-\\textsc{I}{} (Figure~\\ref{fig:single-performance}), and this is reflected in the quality of the learned prior $\\Tilde{p}_\\xi(f)$ in each case (see Appendix~\\ref{app:additional-single}).\nWith respect to the trade-off mentioned in Section~\\ref{sec:prior_covariance}, we find that for small networks, a rich prior over the weights matters more than the full expressiveness of the feature map, making both \\approach-\\textsc{R}{} and \\approach-\\textsc{F}{} appealing.\nHowever, after running further experiments on a deep-learning image-domain problem, this conclusion does not hold for deep networks (see Appendix~\\ref{app:deep}), where keeping an expressive feature map is important (\\approach-\\textsc{I}{} and \\approach-\\textsc{F}{} are appealing in that case).\nThus, \\approach-\\textsc{F}{} is the variant that we retain, for it allows similar or better performances than the other variants in all situations.\n\nNote how \\approach-\\textsc{F}{} outperforms MAML: it achieves much better generalization when decreasing the number of context samples $K$ (Figure~\\ref{fig:single-performance}).\nIndeed, \\approach-\\textsc{F}{} trained with a finite task dataset performs better than MAML with an infinite task dataset: it is able to capture better the common features of the tasks with a smaller task dataset.\n\n\\begin{figure}[t]\n \\centering\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{images\/single_data.pdf}\n \\caption{Examples of context data from in-dist. and OoD tasks}\n \\label{fig:single-data}\n \\end{subfigure}\n \\hfill\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{single_auc.pdf}\n \\caption{AUC for OoD detection}\n \\label{fig:single-auc}\n \\end{subfigure}\n \\hfill\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{single_mse.pdf}\n \\caption{MSE on predictions}\n \\label{fig:single-mse}\n \\end{subfigure}\n \\caption{Unimodal case: Performance of \\textsc{UnLiMiTD}{} for OoD detection and inference, as a function of the number of context datapoints $K$. The training dataset consists of sinusoids, while OoD tasks are lines and quadratic tasks. We compare different variants (\\approach-\\textsc{I}, \\approach-\\textsc{R}{} and \\approach-\\textsc{F}), and against MAML for predictions. We also compare training with a finite and infinite task dataset. Note how \\approach-\\textsc{R}{} and \\approach-\\textsc{F}{} have efficient OoD detection and outperform MAML in predictions. Also, note how MAML trained with an infinite task dataset performs worse than \\approach-\\textsc{R}{} and \\approach-\\textsc{F}{} trained on a finite task dataset.}\n \\label{fig:single-performance}\n\\end{figure}\n\n\\textbf{Multimodal meta-learning: Comparing the mixture model against a single GP.}~\\,\nNext, we consider a multimodal task distribution with training data consisting of sinusoids as well as lines with varying slopes. \nHere, we compare the performance between choosing the mixture structure or the single GP structure (see discussion in Section~\\ref{sec:mixture}): in both cases, we use \\approach-\\textsc{F}.\nMore training details can be found in Appendix~\\ref{app:train-details-multi}).\n\nBoth the OoD detection and the prediction performances are better with the mixture structure than with the single GP structure (Figure~\\ref{fig:multi-performance}), indicating that the mixture model is a useful structure for $\\Tilde{p}_\\xi(f)$.\nThis is reflected in the quality of the learned priors (see Appendix \\ref{app:additional-multi} for qualitative results including samples from the learned priors).\nNote how the single GP structure still performs better than both MAML and MMAML for prediction, especially in the low-data regime.\nThis demonstrates the strength of our probabilistic approach for multimodal meta-learning: even if the probabilistic assumptions are not optimal, the predictions are still accurate and can beat baselines.\n\n\\begin{figure}[t]\n \\centering\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{images\/multi_data.pdf}\n \\caption{Examples of context data from in-dist. and OoD tasks}\n \\label{fig:multi-data}\n \\end{subfigure}\n \\hfill\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{multi_auc.pdf}\n \\caption{AUC for OoD detection}\n \\label{fig:multi-auc}\n \\end{subfigure}\n \\hfill\n \\begin{subfigure}[t]{0.3\\textwidth}\n \\centering\n \\includegraphics[width=\\textwidth]{multi_mse.pdf}\n \\caption{MSE on predictions}\n \\label{fig:multi-mse}\n \\end{subfigure}\n \\caption{Multimodal case: Performance of \\textsc{UnLiMiTD}{} for OoD detection and inference, as a function of the number of context datapoints $K$. The training dataset includes both sines and lines, while OoD tasks are quadratic functions. We compare the different variants (\\approach-\\textsc{F}{} with a single GP or a mixture model), and against MAML\/MMAML for predictions. Note how both versions of \\approach-\\textsc{F}{} yield better predictions than the baselines. In particular, even with a single GP, \\approach-\\textsc{F}{} outperforms the baselines.}\n \\label{fig:multi-performance}\n\\end{figure}\n\n\\section{Conclusion}\\label{sec:conclusion}\n\nWe propose \\textsc{UnLiMiTD}{}, a novel meta-learning algorithm that models the underlying task distribution using a parametric and tuneable distribution, leveraging Bayesian inference with linearized neural networks.\nWe compare three variants, and show that among these, the Fisher-based parameterization, \\approach-\\textsc{F}{}, effectively balances scalability and expressivity, even for deep learning applications.\nWe have demonstrated that (1) our approach makes efficient probabilistic predictions on in-distribution tasks, which compare favorably to, and often outperform, baselines, (2) it allows for effective detection of context data from OoD tasks, and (3) that both these findings continue to hold in the multimodal task-distribution setting.\n\nThere are several avenues for future work. One direction entails understanding how the performance of \\approach-\\textsc{F}{} is impacted if the FIM-based directions are computed too early in the training and the NTK changes significantly afterwards.\nOne could also generalize our approach to non-Gaussian likelihoods, making \\textsc{UnLiMiTD}{} effective for classification tasks.\nFinally, further research can push the limits of multimodal meta-learning, e.g., by implementing non-parametric Bayesian methods to automatically infer an optimal number of clusters, thereby eliminating a hyperparameter of the current approach.\n\n\n\\subsubsection*{Acknowledgements}\nThe authors acknowledge the MIT SuperCloud \\citep{supercloud} and Lincoln Laboratory Supercomputing Center for providing HPC resources that have contributed to the research results reported within this paper.\nThe authors would like to thank MISTI MIT-France for supporting this research.\nC.A. further acknowledges support from Mines Paris Foundation. N.A. acknowledges support from the Edgerton Career Development Professorship.\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Glossary}\\label{glossary}Here we explain some terms that have specific technical definitions in \\gambit.\\begin{description}}\n\\newcommand{\\end{description}}{\\end{description}}\n\n\\newcommand{\\begin{lstlisting}}{\\begin{lstlisting}}\n\\newcommand{\\end{lstlisting}}{\\end{lstlisting}}\n\\newcommand{\\metavarf}[1]{\\textit{\\color{darkgreen}\\footnotesize\\textrm{#1}}}\n\\newcommand{\\metavars}[1]{\\textit{\\color{darkgreen}\\scriptsize\\textrm{#1}}}\n\\newcommand{\\metavarf}{\\metavarf}\n\n\\DeclareMathOperator{\\sign}{sign}\n\\DeclareMathOperator\\erf{erf}\n\n\\newcommand{\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace}{\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace}\n\\newcommand{\\text{M\\eV}\\xspace}{\\text{M\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace}\\xspace}\n\\newcommand{\\text{G\\eV}\\xspace}{\\text{G\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace}\\xspace}\n\\newcommand{\\text{T\\eV}\\xspace}{\\text{T\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace}\\xspace}\n\\newcommand{\\text{pb}\\xspace}{\\text{pb}\\xspace}\n\\newcommand{\\text{fb}\\xspace}{\\text{fb}\\xspace}\n\\newcommand{\\ensuremath{\\pb^{-1}}\\xspace}{\\ensuremath{\\text{pb}\\xspace^{-1}}\\xspace}\n\\newcommand{\\ensuremath{\\fb^{-1}}\\xspace}{\\ensuremath{\\text{fb}\\xspace^{-1}}\\xspace}\n\n\\newcommand{\\ensuremath{p_\\mathrm{T}}\\xspace}{\\ensuremath{p_\\mathrm{T}}\\xspace}\n\\newcommand{\\ensuremath{E_\\mathrm{T}}\\xspace}{\\ensuremath{E_\\mathrm{T}}\\xspace}\n\\newcommand{\\ensuremath{E_\\mathrm{T}^\\mathrm{\\mspace{1.5mu}miss}}\\xspace}{\\ensuremath{E_\\mathrm{T}^\\mathrm{\\mspace{1.5mu}miss}}\\xspace}\n\\newcommand{\\etmissx}[1]{\\ensuremath{E_\\mathrm{T}^\\mathrm{\\mspace{1.5mu}miss,{#1}}}\\xspace}\n\\newcommand{\\ensuremath{H_\\mathrm{T}}\\xspace}{\\ensuremath{H_\\mathrm{T}}\\xspace}\n\\newcommand{\\ensuremath{\\Delta\\phi}\\xspace}{\\ensuremath{\\Delta\\phi}\\xspace}\n\\newcommand{\\scriptscriptstyle}{\\scriptscriptstyle}\n\\newcommand{m_{\\sss S}}{m_{\\scriptscriptstyle S}}\n\\newcommand{\\overline{m}_{\\sss S}}{\\overline{m}_{\\scriptscriptstyle S}}\n\\newcommand{\\lambda_{h\\sss S}}{\\lambda_{h\\scriptscriptstyle S}}\n\\newcommand{\\lambda_{\\sss S}}{\\lambda_{\\scriptscriptstyle S}}\n\\newcommand{\\lambda_{h}}{\\lambda_{h}}\n\\newcommand{m_h}{m_h}\n\\newcommand{\\overline{m}_h}{\\overline{m}_h}\n\\newcommand{$\\overline{DR}$\\xspace}{$\\overline{DR}$\\xspace}\n\\newcommand{\\DR}{$\\overline{DR}$\\xspace}\n\\newcommand{$\\overline{MS}$\\xspace}{$\\overline{MS}$\\xspace}\n\\newcommand{\\overline{MS}}{\\overline{MS}}\n\\newcommand{\\text{CL}\\xspace}{\\text{CL}\\xspace}\n\\newcommand{\\ensuremath{\\CL_{s}}\\xspace}{\\ensuremath{\\text{CL}\\xspace_{s}}\\xspace}\n\\newcommand{\\ensuremath{\\CL_{s+b}}\\xspace}{\\ensuremath{\\text{CL}\\xspace_{s+b}}\\xspace}\n\\newcommand{\\ensuremath{\\mathrm{BR}}\\xspace}{\\ensuremath{\\mathrm{BR}}\\xspace}\n\\newcommand{$\\mathbb{Z}_2$\\xspace}{$\\mathbb{Z}_2$\\xspace}\n\\newcommand{$\\mathbb{Z}_3$\\xspace}{$\\mathbb{Z}_3$\\xspace}\n\n\\newcommand{i.e.\\ }{i.e.\\ }\n\\newcommand{e.g.\\ }{e.g.\\ }\n\\newcommand{ATLAS\\xspace}{ATLAS\\xspace}\n\\newcommand{CMS\\xspace}{CMS\\xspace}\n\\newcommand{\\textsf{GAMBIT}\\xspace}{\\textsf{GAMBIT}\\xspace}\n\\newcommand{1.0.0}{1.0.0}\n\\newcommand{\\gambit \\textsf{\\gambitversion}\\xspace}{\\textsf{GAMBIT}\\xspace \\textsf{1.0.0}\\xspace}\n\\newcommand{\\textsf{DarkBit}\\xspace}{\\textsf{DarkBit}\\xspace}\n\\newcommand{\\textsf{ColliderBit}\\xspace}{\\textsf{ColliderBit}\\xspace}\n\\newcommand{\\textsf{FlavBit}\\xspace}{\\textsf{FlavBit}\\xspace}\n\\newcommand{\\textsf{SpecBit}\\xspace}{\\textsf{SpecBit}\\xspace}\n\\newcommand{\\textsf{DecayBit}\\xspace}{\\textsf{DecayBit}\\xspace}\n\\newcommand{\\textsf{PrecisionBit}\\xspace}{\\textsf{PrecisionBit}\\xspace}\n\\newcommand{\\textsf{ScannerBit}\\xspace}{\\textsf{ScannerBit}\\xspace}\n\\newcommand{\\textsf{ExampleBit\\_A}\\xspace}{\\textsf{ExampleBit\\_A}\\xspace}\n\\newcommand{\\textsf{ExampleBit\\_B}\\xspace}{\\textsf{ExampleBit\\_B}\\xspace}\n\\newcommand{\\textsf{NeutrinoBit}\\xspace}{\\textsf{NeutrinoBit}\\xspace}\n\\newcommand{\\textsf{BOSS}\\xspace}{\\textsf{BOSS}\\xspace}\n\\newcommand{\\gambit}{\\textsf{GAMBIT}\\xspace}\n\\newcommand{\\darkbit}{\\textsf{DarkBit}\\xspace}\n\\newcommand{\\textsf{OpenMP}\\xspace}{\\textsf{OpenMP}\\xspace}\n\\newcommand{\\textsf{MPI}\\xspace}{\\textsf{MPI}\\xspace}\n\\newcommand{\\textsf{POSIX}\\xspace}{\\textsf{POSIX}\\xspace}\n\\newcommand{\\textsf{BuckFast}\\xspace}{\\textsf{BuckFast}\\xspace}\n\\newcommand{\\textsf{Delphes}\\xspace}{\\textsf{Delphes}\\xspace}\n\\newcommand{\\textsf{EOS}\\xspace}{\\textsf{EOS}\\xspace}\n\\newcommand{\\textsf{Flavio}\\xspace}{\\textsf{Flavio}\\xspace}\n\\newcommand{\\textsf{Pythia}\\xspace}{\\textsf{Pythia}\\xspace}\n\\newcommand{\\textsf{Pythia\\,8}\\xspace}{\\textsf{Pythia\\,8}\\xspace}\n\\newcommand{\\textsf{PythiaEM}\\xspace}{\\textsf{PythiaEM}\\xspace}\n\\newcommand{\\textsf{Prospino}\\xspace}{\\textsf{Prospino}\\xspace}\n\\newcommand{\\textsf{NLL-fast}\\xspace}{\\textsf{NLL-fast}\\xspace}\n\\newcommand{\\textsf{MadGraph}\\xspace}{\\textsf{MadGraph}\\xspace}\n\\newcommand{\\textsf{MadGraph5\\_aMC@NLO}\\xspace}{\\textsf{MadGraph5\\_aMC@NLO}\\xspace}\n\\newcommand{\\textsf{FastJet}\\xspace}{\\textsf{FastJet}\\xspace}\n\\newcommand{\\textsf{SModelS}\\xspace}{\\textsf{SModelS}\\xspace}\n\\newcommand{\\textsf{FastLim}\\xspace}{\\textsf{FastLim}\\xspace}\n\\newcommand{\\textsf{CheckMATE}\\xspace}{\\textsf{CheckMATE}\\xspace}\n\\newcommand{\\textsf{HiggsBounds}\\xspace}{\\textsf{HiggsBounds}\\xspace}\n\\newcommand{\\textsf{SUSYPope}\\xspace}{\\textsf{SUSYPope}\\xspace}\n\\newcommand{\\textsf{HiggsSignals}\\xspace}{\\textsf{HiggsSignals}\\xspace}\n\\newcommand{\\textsf{DarkSUSY}\\xspace}{\\textsf{DarkSUSY}\\xspace}\n\\newcommand{\\ds}{\\textsf{DarkSUSY}\\xspace}\n\\newcommand{\\textsf{WimpSim}\\xspace}{\\textsf{WimpSim}\\xspace}\n\\newcommand{\\textsf{3-BIT-HIT}\\xspace}{\\textsf{3-BIT-HIT}\\xspace}\n\\newcommand{\\textsf{PPPC4DMID}\\xspace}{\\textsf{PPPC4DMID}\\xspace}\n\\newcommand{\\textsf{micrOMEGAs}\\xspace}{\\textsf{micrOMEGAs}\\xspace}\n\\newcommand{\\textsf{micrOMEGAs}\\xspace}{\\textsf{micrOMEGAs}\\xspace}\n\\newcommand{\\textsf{Rivet}\\xspace}{\\textsf{Rivet}\\xspace}\n\\newcommand{\\textsf{Feynrules}\\xspace}{\\textsf{Feynrules}\\xspace}\n\\newcommand{\\textsf{FeynHiggs}\\xspace}{\\textsf{FeynHiggs}\\xspace}\n\\newcommand{\\feynhiggs}{\\textsf{FeynHiggs}\\xspace}\n\\newcommand{\\textsf{EOS}\\xspace}{\\textsf{EOS}\\xspace}\n\\newcommand{\\textsf{FlavorKit}\\xspace}{\\textsf{FlavorKit}\\xspace}\n\\newcommand\\FS{\\FlexibleSUSY}\n\\newcommand\\flexiblesusy{\\FlexibleSUSY}\n\\newcommand\\FlexibleSUSY{\\textsf{FlexibleSUSY}\\xspace}\n\\newcommand\\FlexibleEFTHiggs{\\textsf{FlexibleEFTHiggs}\\xspace}\n\\newcommand\\SOFTSUSY{\\textsf{SOFTSUSY}\\xspace}\n\\newcommand\\SUSPECT{\\textsf{SuSpect}\\xspace}\n\\newcommand\\NMSSMCalc{\\textsf{NMSSMCALC}\\xspace}\n\\newcommand\\NMSSMTools{\\textsf{NMSSMTools}\\xspace}\n\\newcommand\\NMSPEC{\\textsf{NMSPEC}\\xspace}\n\\newcommand\\NMHDECAY{\\textsf{NMHDECAY}\\xspace}\n\\newcommand\\HDECAY{\\textsf{HDECAY}\\xspace}\n\\newcommand\\prophecy{\\textsf{PROPHECY4F}\\xspace}\n\\newcommand\\SDECAY{\\textsf{SDECAY}\\xspace}\n\\newcommand\\SUSYHIT{\\textsf{SUSY-HIT}\\xspace}\n\\newcommand\\susyhd{\\textsf{SUSYHD}\\xspace}\n\\newcommand\\HSSUSY{\\textsf{HSSUSY}\\xspace}\n\\newcommand\\susyhit{\\SUSYHIT}\n\\newcommand\\gmtwocalc{\\textsf{GM2Calc}\\xspace}\n\\newcommand\\SARAH{\\textsf{SARAH}\\xspace}\n\\newcommand\\SPheno{\\textsf{SPheno}\\xspace}\n\\newcommand\\spheno{\\SPheno}\n\\newcommand\\superiso{\\textsf{SuperIso}\\xspace}\n\\newcommand\\superisofour{\\textsf{SuperIso 4}\\xspace}\n\\newcommand\\heplike{\\textsf{HEPLike}\\xspace}\n\\newcommand\\SFOLD{\\textsf{SFOLD}\\xspace}\n\\newcommand\\HFOLD{\\textsf{HFOLD}\\xspace}\n\\newcommand\\FeynHiggs{\\textsf{FeynHiggs}\\xspace}\n\\newcommand\\Mathematica{\\textsf{Mathematica}\\xspace}\n\\newcommand\\Kernel{\\textsf{Kernel}\\xspace}\n\\newcommand\\WSTP{\\textsf{WSTP}\\xspace}\n\\newcommand\\lilith{\\textsf{Lilith}\\xspace}\n\\newcommand\\nulike{\\textsf{nulike}\\xspace}\n\\newcommand\\gamLike{\\textsf{gamLike}\\xspace}\n\\newcommand\\gamlike{\\gamLike}\n\\newcommand\\daFunk{\\textsf{daFunk}\\xspace}\n\\newcommand\\pippi{\\textsf{pippi}\\xspace}\n\\newcommand\\MultiNest{\\textsf{MultiNest}\\xspace}\n\\newcommand\\multinest{\\MultiNest}\n\\newcommand\\great{\\textsf{GreAT}\\xspace}\n\\newcommand\\twalk{\\textsf{T-Walk}\\xspace}\n\\newcommand\\diver{\\textsf{Diver}\\xspace}\n\\newcommand\\ddcalc{\\textsf{DDCalc}\\xspace}\n\\newcommand\\tpcmc{\\textsf{TPCMC}\\xspace}\n\\newcommand\\nest{\\textsf{NEST}\\xspace}\n\\newcommand\\luxcalc{\\textsf{LUXCalc}\\xspace}\n\\newcommand\\xx{\\raisebox{0.2ex}{\\smaller ++}\\xspace}\n\\newcommand\\Cpp{\\textsf{C\\xx}\\xspace}\n\\newcommand\\Cppeleven{\\textsf{C\\raisebox{0.2ex}{\\smaller ++}11}\\xspace}\n\\newcommand\\plainC{\\textsf{C}\\xspace}\n\\newcommand\\Python{\\textsf{Python}\\xspace}\n\\newcommand\\python{\\Python}\n\\newcommand\\Fortran{\\textsf{Fortran}\\xspace}\n\\newcommand\\YAML{\\textsf{YAML}\\xspace}\n\\newcommand\\Yaml{\\YAML}\n\\newcommand{\\textsf{Capt'n General}\\xspace}{\\textsf{Capt'n General}\\xspace}\n\n\\newcommand\\beq{\\begin{equation}}\n\\newcommand\\eeq{\\end{equation}}\n\n\\newcommand{\\mail}[1]{\\href{mailto:#1}{#1}}\n\\renewcommand{\\url}[1]{\\href{#1}{#1}}\n\n\\newcommand{\\TODO}[1]{\\textbf{\\textcolor{red}{To do: #1}}}\n\\newcommand{\\Pat}[1]{\\textbf{\\color{teal}Pat: #1}}\n\\newcommand{\\Anders}[1]{\\textbf{\\color{brown}Anders: #1}}\n\\newcommand{\\Martin}[1]{\\textbf{\\color{orange}Martin: #1}}\n\n\n\\def\\alpha_t{\\alpha_t}\n\\def\\alpha_b{\\alpha_b}\n\\def\\alpha_s{\\alpha_s}\n\\def\\alpha_{\\tau}{\\alpha_{\\tau}}\n\\def\\mathcal{O}(\\at){\\mathcal{O}(\\alpha_t)}\n\\def\\mathcal{O}(\\ab){\\mathcal{O}(\\alpha_b)}\n\\def\\mathcal{O}(\\atau){\\mathcal{O}(\\alpha_{\\tau})}\n\\def\\mathcal{O}(\\at\\ab){\\mathcal{O}(\\alpha_t\\alpha_b)}\n\\def\\mathcal{O}(\\at\\as){\\mathcal{O}(\\alpha_t\\alpha_s)}\n\\def\\mathcal{O}(\\ab\\as){\\mathcal{O}(\\alpha_b\\alpha_s)}\n\\def\\mathcal{O}(\\at\\ab + \\ab^2){\\mathcal{O}(\\alpha_t\\alpha_b + \\alpha_b^2)}\n\\def\\mathcal{O}(\\at^2 + \\at\\ab + \\ab^2){\\mathcal{O}(\\alpha_t^2 + \\alpha_t\\alpha_b + \\alpha_b^2)}\n\\def\\mathcal{O}(\\at\\as + \\at^2){\\mathcal{O}(\\alpha_t\\alpha_s + \\alpha_t^2)}\n\\def\\mathcal{O}(\\at\\as +\\ab\\as){\\mathcal{O}(\\alpha_t\\alpha_s +\\alpha_b\\alpha_s)}\n\\def\\mathcal{O}(\\at\\as + \\at^2 +\\ab\\as){\\mathcal{O}(\\alpha_t\\alpha_s + \\alpha_t^2 +\\alpha_b\\alpha_s)}\n\\def\\mathcal{O}(\\at^2){\\mathcal{O}(\\alpha_t^2)}\n\\def\\mathcal{O}(\\ab^2){\\mathcal{O}(\\alpha_b^2)}\n\\def\\mathcal{O}(\\atau^2){\\mathcal{O}(\\alpha_{\\tau}^2)}\n\\def\\mathcal{O}(\\ab \\atau){\\mathcal{O}(\\alpha_b \\alpha_{\\tau})}\n\\def\\mathcal{O}((\\at+\\ab)^2){\\mathcal{O}((\\alpha_t+\\alpha_b)^2)}\n\\def\\mathcal{O}(\\as){\\mathcal{O}(\\alpha_s)}\n\\def\\mathcal{O}(\\atau^2 +\\ab \\atau ){\\mathcal{O}(\\alpha_{\\tau}^2 +\\alpha_b \\alpha_{\\tau} )}\n\n\n\\begin{document}\n\n\\begin{frontmatter}\n\n\n\n\\title{GAMBIT and its Application in the Search for Physics Beyond the Standard Model}\n\n\n\\author[imperial,oslo]{Anders Kvellestad}\n\\author[imperial,uq]{Pat Scott}\n\\author[adelaide]{Martin White}\n\\address[imperial]{Department of Physics, Imperial College London, Blackett Laboratory, Prince Consort Road, London SW7 2AZ, UK}\n\\address[oslo]{Department of Physics, University of Oslo, N-0316 Oslo, Norway}\n\\address[uq]{School of Mathematics and Physics, The University of Queensland, St.\\ Lucia, Brisbane, QLD 4072, Australia}\n\\address[adelaide]{Department of Physics, University of Adelaide, Adelaide, SA 5005, Australia}\n\n\\begin{abstract}\nThe Global and Modular Beyond-Standard Model Inference Tool (\\textsf{GAMBIT}\\xspace) is an open source software framework for performing global statistical fits of particle physics models, using a wide range of particle and astroparticle data. In this review, we describe the design principles of the package, the statistical and sampling frameworks, the experimental data included, and the first two years of physics results generated with it. This includes supersymmetric models, axion theories, Higgs portal dark matter scenarios and an extension of the Standard Model to include right-handed neutrinos. Owing to the broad spectrum of physics scenarios tackled by the \\textsf{GAMBIT}\\xspace community, this also serves as a convenient, self-contained review of the current experimental and theoretical status of the most popular models of dark matter.\n\\end{abstract}\n\n\n\\begin{keyword}\n Supersymmetry \\sep Axions \\sep Global fits \\sep Right-handed neutrinos \\sep Higgs portal dark matter\n\n\n\n\\end{keyword}\n\n\\end{frontmatter}\n\n\n\\section{Introduction}\n\nThe core of the scientific method in the physical sciences is the identification of mathematical theories that describe some aspect of our Universe in terms of a number of free parameters. The use of global statistical fits, in either a Bayesian or frequentist framework, allows us to find the preferred parameter values of a candidate theory given experimental data, and to compare the abilities of different models to describe that data. Our current knowledge of particle physics is enshrined in the Standard Model (SM), and global fit techniques are routinely used to provide the most accurate estimates the parameters of the neutrino sector~\\cite{Bari13,Tortola14,NuFit15}, the CKM matrix~\\cite{CKMFitter}, and the electroweak sector~\\cite{ZFitter,GFitter11}.\n\nDespite its incredible successes in explaining experimental data, the SM still faces a number of experimental and theoretical challenges. Many if not all of these can be explained by new physics Beyond the Standard Model (BSM). Such physics could show up in a number of experiments, including direct searches for new particles at high energy particle colliders~\\cite{ATLAS_diphoton,ATLAS15,CMS_SMS}, measurements of rare Standard Model processes~\\cite{gm2exp,BelleII,CMSLHCb_Bs0mumu}, direct searches for dark matter~\\cite{XENON2013,PICO60,LUX2016}, indirect astroparticle searches for distant annihilation or decay of dark matter~\\cite{BringmannWeniger,LATdwarfP8,IC79_SUSY}, and cosmological observations~\\cite{Planck15cosmo,Slatyer15a,keVsterile_whitepaper}. Unfortunately, despite the existence of many candidate theories beyond the SM, there is no unambiguous prediction of what we expect to observe, or in which experimental field we expect to observe it. It is therefore highly likely that the next theory of particle physics will have to be pieced together by combining clues from a number of disparate fields and experiments. In the process, it is essential to also consistently combine \\emph{null} results in experiments that had the potential to discover a given candidate theory, but failed to do so. Even in the complete absence of positive discoveries in the near future, it is essential to determine which candidate BSM theories are now comfortably excluded, and which regions of which candidate theories are now the most amenable to future discovery.\n\nGlobal fits of BSM theories have thus been a very active area of research for well over a decade \\cite{Baltz04,Allanach06,SFitter, Ruiz06}, with increases in computing power opening the option of exploring models with larger and larger parameter spaces. Nevertheless, it remains a considerable challenge to efficiently explore the high-dimensional parameter spaces of candidate theories whilst rigorously calculating likelihoods for a large range of experiments, each of which may require a costly simulation procedure. To further complicate matters, one must consistently handle systematic uncertainties that may be correlated across different datasets, resulting from either instrumental effects, or our imprecise knowledge of the nuclear, astro- or particle physics relevant to a given set of experiments. Prior to 2017, most global fits were focussed on supersymmetric theories, involving dedicated software that was built from the ground up with a knowledge of the supersymmetric parameters~\\cite{Baltz04,Allanach06,SFitter, Ruiz06,Strege15,Fittinocoverage,Catalan:2015cna,MasterCodeMSSM10,2007NewAR..51..316T,2007JHEP...07..075R,Roszkowski09a,Martinez09,Roszkowski09b,Roszkowski10,Scott09c,BertoneLHCDD,SBCoverage,Nightmare,BertoneLHCID,IC22Methods,SuperbayesXENON100,SuperBayesGC, Buchmueller08,Buchmueller09,MasterCodemSUGRA,MasterCode11,MastercodeXENON100,MastercodeHiggs,Buchmueller:2014yva,Bagnaschi:2016afc,Bagnaschi:2016xfg,Allanach:2007qk,Abdussalam09a,Abdussalam09b,Allanach11b,Allanach11a,Farmer13,arXiv:1212.4821,Fowlie13,Henrot14,Kim:2013uxa,arXiv:1503.08219,arXiv:1604.02102,Han:2016gvr, Bechtle:2014yna, arXiv:1405.4289, arXiv:1402.5419, MastercodeCMSSM, arXiv:1312.5233, arXiv:1310.3045, arXiv:1309.6958, arXiv:1307.3383, arXiv:1304.5526, arXiv:1212.2886, Strege13, Gladyshev:2012xq, Kowalska:2012gs, Mastercode12b, arXiv:1207.1839, arXiv:1207.4846, Roszkowski12, SuperbayesHiggs, Fittino12, Mastercode12, arXiv:1111.6098, Fittino, Trotta08, Fittino06,\narXiv:1608.02489, arXiv:1507.07008, Mastercode15, arXiv:1506.02499, arXiv:1504.03260, Mastercode17}. These results typically covered low-dimensional subsets of the minimal supersymmetric SM (MSSM) or, in some cases the next-to-minimal variant, with relatively few global studies of other theories completed~\\cite{Cheung:2012xb,Arhrib:2013ela,Sming14,Chowdhury15,Liem16,LikeDM,Banerjee:2016hsk,Matsumoto:2016hbs,Cuoco:2016jqt,Cacchio:2016qyh,BertoneUED,Chiang:2018cgb,hepfit,Matsumoto:2018acr}.\n\nIn 2017, the \\textsf{GAMBIT}\\xspace collaboration released the Global and Modular Beyond-Standard Model Inference Tool (\\textsf{GAMBIT}\\xspace) \\cite{gambit}, an open-source package able to produce results in both the Bayesian and frequentist statistical frameworks, and easily extendible to new BSM models and new experimental datasets. A fully modular design enables much of the code to be reused when changing the theoretical model of interest. \\textsf{GAMBIT}\\xspace includes a wide variety of efficient sampling algorithms for posterior evaluation and optimisation, and ensures computational efficiency through massive, multi-level parallelisation, both of the sampling algorithms and individual likelihood calculations.\n\nThe purpose of this article is to give a brief introduction to the \\textsf{GAMBIT}\\xspace software and science programme, reviewing the most important results obtained with \\textsf{GAMBIT}\\xspace in the first few years since its initial release. These serve to illustrate the versatility of the code in attacking completely different BSM models, and the constraining power of the highly detailed and rigorous simulations of different particle and astroparticle datasets in \\textsf{GAMBIT}\\xspace. Given the centrality of dark matter (DM) in the current search for BSM physics, this review also serves as a convenient summary of the status of the most widely-studied DM candidates.\n\nIn Section~\\ref{sec:gambit}, we describe the structure and design of the \\textsf{GAMBIT}\\xspace package, including the core framework, the means by which \\textsf{GAMBIT}\\xspace supports generic BSM models, the sampling and statistics module, and the various physics modules able to produce theoretical predictions and experimental likelihoods. In Section~\\ref{sec:physics}, we summarise the results of recent \\textsf{GAMBIT}\\xspace global fits of various supersymmetric theories, Higgs portal and axion DM models, and a right-handed neutrino extension of the SM. We then conclude in Section~\\ref{sec:summary}.\n\n\\section{The \\textsf{GAMBIT}\\xspace software}\n\\label{sec:gambit}\n\nThe core \\textsf{GAMBIT}\\xspace software is written in \\Cppeleven, but interfaces seamlessly with extensions and existing physics codes written in \\python, \\Mathematica, \\Fortran and \\plainC. Since the release of \\textsf{GAMBIT}\\xspace \\textsf{1.0.0} in 2017 \\cite{gambit}, the most notable updates have been versions \\textsf{1.1} (adding support for \\Mathematica) \\cite{gambit_addendum}, \\textsf{1.2} (adding support for \\python and higher-spin Higgs portal models) \\cite{HP}, \\textsf{1.3} (adding support for axion and axion-like particles) \\cite{Axions} and \\textsf{1.4} (adding support for right-handed neutrinos) \\cite{RHN}. The current public release is \\textsf{v1.4.2}. The source code is openly available from \\href{http:\/\/gambit.hepforge.org}{http:\/\/gambit.hepforge.org} under the 3-clause BSD license.\n\n\\subsection{Core design}\n\\label{sec:core}\n\nThe core principles of \\textsf{GAMBIT}\\xspace's software design are modularity and flexibility. All theoretical predictions and experimental likelihood evaluations are separated into a series of smaller, self-contained sub-calculations, with each sub-calculation represented by a single function. Each function is assigned a metadata string that identifies the physical quantity that the function is able to calculate. Examples might be the mass of the lightest Higgs boson, or the likelihood for the latest run of the LUX direct detection experiment. Functions are further tagged with additional metadata strings indicating any other physical inputs required for them to run. In the case of the Higgs mass, one might require e.g. the SM electroweak vacuum expectation value and the masses and couplings of various other particles. In the case of the LUX likelihood, one might require the number of events observed by LUX, its detector efficiency as a function of nuclear recoil energy, and the theoretically predicted event rate.\nAt runtime, \\textsf{GAMBIT}\\xspace identifies which functions are actually required for the analysis of a given theory, and connects them dynamically in order to enable the calculation in the most efficient manner possible.\n\nThe individual functions are grouped together according to physics theme, into seven different \\textbf{physics modules}. We describe these specific modules in Section \\ref{sec:modules} below. Individual functions are thus referred to as \\textbf{module functions}. The module functions are the true building blocks of a \\textsf{GAMBIT}\\xspace analysis, allowing the code to automatically adapt itself to incorporate new observables, likelihoods, theories and experimental datasets. The metadata string associated with the output of a module function is referred to as its \\textbf{capability}, and the metadata associated with the required inputs are referred to as \\textbf{dependencies}. The process of dynamically connecting the outputs of module functions to the inputs of others at runtime thus consists of matching dependency strings to the capability strings of other functions (and ensuring that their \\Cpp types also match). This process is known as \\textbf{dependency resolution}, and is performed by the \\textsf{GAMBIT}\\xspace dependency resolver.\n\n\\begin{figure}[t]\n\t\\centering\n\t\\includegraphics[width = 0.75\\textwidth]{figures\/dep_tree}\n\t\\caption{An example \\textsf{GAMBIT}\\xspace dependency tree for a simple fit of flavour Wilson coefficients to $b\\to s\\gamma$ and $B\\to ll$ data. Boxes (graph nodes) correspond to single module functions. Function capabilities are marked in red, and return types of the functions, their actual function names and enveloping modules are indicated in black. Arrows (graph edges) indicate the direction of information flow, from the capability (output) of one function to the dependencies (inputs) of others. The input file used to instigate this fit (\\texttt{WC\\_lite.yaml}) is one of the example files distributed with \\textsf{GAMBIT}\\xspace. This particular fit makes use of the \\textsf{GAMBIT}\\xspace modules \\textsf{FlavBit}\\xspace, \\textsf{SpecBit}\\xspace and \\textsf{DecayBit}\\xspace, as well as the backend (external package) \\superiso \\textsf{4.1} \\cite{Mahmoudi:2007vz,Mahmoudi:2008tp,Mahmoudi:2009zz}.}\n\t\\label{fig:dep_tree}\n\\end{figure}\n\nThe result is a (potentially extremely large) directed graph connecting the outputs and inputs of different module functions, known as a \\textbf{dependency tree}. The module functions themselves constitute graph nodes, and their resolved dependencies graph edges. An example of such a graph is shown in Fig.\\ \\ref{fig:dep_tree}. For such a graph to constitute a viable computational pathway to all theoretical predictions and experimental likelihoods of interest, two things are required. The first is for the actual numerical values of some dependencies to be known in advance. These are the parameters of the theory that the user wishes to analyse, and must be chosen `from on high' before the dependency tree can be evaluated. These are selected by the user's chosen statistical parameter sampling algorithm, discussed below in Section \\ref{sec:stats}.\n\nFrom the values of a model's parameters, all other intermediate quantities can be obtained, as long as the second criterion is also met. This condition is that no closed loops exist in the graph, i.e.\\ there are no dependencies of any module function upon things that can only be computed by knowing the result of the function. Many algorithms exist within graph theory for taking such a directed acyclic graph and obtaining a linear ordering of its nodes that respects the underlying structure of the graph. \\textsf{GAMBIT}\\xspace uses the \\textsf{Boost::Graph} library to obtain such an ordering, and then employs that ordering to evaluate the module functions in turn. This ensures that all module functions run before any other functions that depend upon their results. Within topologically equivalent subsets of the ordering, \\textsf{GAMBIT}\\xspace also further dynamically optimises the module function evaluation order for speed, according to previous function evaluation times and likelihoods.\n\nModule functions may also make use of functions provided by external packages, or \\textbf{backends}. These are also connected dynamically at runtime to module functions by the dependency resolver, in much the same way as it ensures that dependencies upon the results of other module functions are fulfilled. This layer of abstraction allows \\textsf{GAMBIT}\\xspace to provide its module functions with seamless and interchangeable access to functions from external codes written in \\plainC, \\Cpp, \\python\\textsf{2}, \\python\\textsf{3}, \\Mathematica and all variants of \\Fortran. The \\textsf{GAMBIT}\\xspace build system allows users to select and automatically download, configure and build whatever combination of backends they prefer to use, and the dependency resolver automatically adapts to the presence or absence of different backends when selecting which functions to connect to others. Backends presently supported in version \\textsf{1.4.2} of \\textsf{GAMBIT}\\xspace are \\textsf{Capt'n General}\\xspace \\cite{HP}, \\ds \\cite{darksusy4,darksusy}, \\ddcalc \\cite{DarkBit,HP}, \\textsf{FeynHiggs}\\xspace \\cite{Heinemeyer:1998yj,Heinemeyer:1998np,Degrassi:2002fi,Frank:2006yh,Hahn:2013ria}, \\gamlike \\cite{DarkBit}, \\gmtwocalc \\cite{gm2calc}, \\textsf{HiggsBounds}\\xspace \\cite{Bechtle:2008jh,Bechtle:2011sb,Bechtle:2013wla}, \\textsf{HiggsSignals}\\xspace \\cite{HiggsSignals}, \\textsf{micrOMEGAs}\\xspace \\cite{Belanger:2001fz,Belanger:2004yn,Belanger:2006is,Belanger:2008sj,Belanger:2010gh,Belanger:2013oya,micromegas}, \\nulike \\cite{IC22Methods,IC79_SUSY}, \\textsf{Pythia}\\xspace \\cite{Sjostrand:2006za,Sjostrand:2014zea}, \\spheno \\cite{Porod:2003um,Porod:2011nf}, \\superiso \\cite{Mahmoudi:2007vz,Mahmoudi:2008tp,Mahmoudi:2009zz}, \\susyhd \\cite{Vega:2015fna} and \\susyhit \\cite{Djouadi:2006bz}. Many more are also already supported in the current development version, which will be released in 2020.\n\n\\begin{figure}[tbp]\n\t\\centering\n\t\\includegraphics[width = \\textwidth]{figures\/BasicStructure}\n\t\\caption{The overall structure of a \\textsf{GAMBIT}\\xspace run, illustrating the roles of the input \\YAML file, modules, module functions, capabilities, dependencies, backends, backend functions, dependency resolver, hierarchical model database and the sampling machinery. The user specifies one or more models to scan in the input \\YAML file, and chooses likelihoods and observables to compute in the scan, making their choice by capability rather than by choosing specific functions. The dependency resolver automatically identifies and connects appropriate module and backend functions in order to facilitate the computation of the requested likelihoods and observables, and the scanning machinery (\\textsf{ScannerBit}\\xspace) selects parameter combinations to pass through the resulting dependency tree. From \\cite{gambit}.}\n\t\\label{fig:corechain}\n\\end{figure}\n\nActual runs of \\textsf{GAMBIT}\\xspace are driven by a single input file, in \\YAML format. In this file, the user selects the model(s) to analyse, gives details of which algorithms to use in order to sample the models' parameters, and provides a list of all likelihoods and physical observables that should be calculated in the scan. The model parameter values constitute one boundary condition for dependency resolution (the dependency tree must begin from the parameters), and the target likelihoods and observables the other (the final outputs of the tree must be the required likelihoods and observables). The dependency resolver is then responsible for identifying and filling in all the required steps in between. To help direct this process and break degeneracies in the valid choices available to the dependency resolver at each step, the \\YAML input file may also set rules that the dependency resolver must respect. These may be e.g.\\ restrictions about which functions should be selected to fill which specific dependencies, or which version of a given backend should be used throughout the run. These rules can be arbitrarily complicated, general or specific. A rule can also contain explicit keyword options that will be passed to all module functions that fulfil the rule, allowing enormous control to be exercised over the details of the individual calculations from a single input file.\n\nThe core design elements of \\textsf{GAMBIT}\\xspace described so far in this section are module functions, backend abstraction, dependency resolution, and an input format that borders on its own programming language. Together, these combine to provide an extremely flexible and extendible framework for performing global analyses of theories for BSM physics. Fig.\\ \\ref{fig:corechain} illustrates how all of these features work together to enable a \\textsf{GAMBIT}\\xspace scan. Further technical details can be found in Ref.\\ \\cite{gambit}.\n\n\n\\subsection{Model support}\n\\label{sec:models}\n\nAnother feature illustrated in Fig.\\ \\ref{fig:corechain} is the \\textsf{GAMBIT}\\xspace hierarchical model database. Models are defined both in terms of their parameters, and in terms of their relationships to each other via parameter translation routines. Models may descend from one another, meaning that a parameter combination in a child model can be translated `up' its family tree to a point in an appropriate subspace of its parent model, or in any other more distant ancestor model. Cross-family `friend' translation pathways can also be defined. These pathways allow module functions to be designed to work with one model, but to be used with another model without further alteration, so long as a translation pathway exists from one model to the other.\n\nModule functions, backend functions and all rules set in \\YAML files can be endowed with model-specific restrictions. This allows the model-dependence of every sub-calculation to be tracked explicitly, and for the dependency resolver to explicitly ensure that the entire dependency tree of every scan is validated for use with the model under investigation.\n\nThe complete model database of \\textsf{GAMBIT}\\xspace \\textsf{1.4.2} is shown in Fig.\\ \\ref{fig:model_tree}.\n\n\\begin{sidewaysfigure}[tbp]\n\t\\centering\n\t\\includegraphics[width = \\textheight]{figures\/model_tree}\\vspace{2mm}\n\t\\caption{Hierarchical model database of \\textsf{GAMBIT}\\xspace \\textsf{1.4.2}. Models are shown as boxes, child-to-parent translations as black arrows, and friend translations as red arrows.}\n\t\\label{fig:model_tree}\n\\end{sidewaysfigure}\n\n\n\\subsection{Sampling and statistics}\n\\label{sec:stats}\n\nIn carrying out a global statistical analysis of a BSM theory, one may be interested in determining which parameter combinations are able to explain the totality of observed data within a given model, and to what extent -- or one may be more interested in using the experimental data to choose between entire theories. The first of these tasks is parameter estimation, whereas the second is model comparison. There are two philosophically distinct ways of posing both these questions:\n\\begin{enumerate}\n\\item How probable is it that we would have observed the data that we have, if a model and a specific combination of its parameters were true?\n\\item How probable is it that a model (or a specific combination of its parameters) is true, given the data that we have observed to date?\n\\end{enumerate}\nQuestion 1 concerns frequentist statistics, whereas Question 2 is fundamentally Bayesian.\n\nIn the context of parameter estimation, the choice of question dictates whether the appropriate quantity to consider is the frequentist profile likelihood, or the Bayesian posterior. The profile likelihood for some parameters of interest $\\boldsymbol{\\theta}$ is the maximum value of the likelihood at each parameter combination $\\boldsymbol{\\theta}$, regardless of the values of any other parameters $\\boldsymbol{\\alpha}$:\n\\begin{equation}\n \\hat{\\mathcal{L}}(\\boldsymbol{\\theta}) = \\mathrm{max}_{\\boldsymbol{\\alpha}}\\,\\mathcal{L}(\\boldsymbol{\\theta},\\boldsymbol{\\alpha}),\n\\end{equation}\nwhere the parameters $\\boldsymbol{\\alpha}$ are other `nuisance' parameters not of direct interest. Conversely, the Bayesian posterior probability distribution for the parameters $\\boldsymbol{\\theta}$ is given by Bayes' Theorem as the integral of the likelihood over $\\boldsymbol\\alpha$, weighted by ones prior belief $\\pi(\\boldsymbol{\\theta},\\boldsymbol{\\alpha})$ as to the plausibility of different values of $\\boldsymbol{\\theta}$ and $\\boldsymbol{\\alpha}$:\n\\begin{equation}\n\\mathcal{P}(\\boldsymbol\\theta) = \\int \\mathcal{P}(\\boldsymbol\\theta,\\boldsymbol{\\alpha})\\,\\mathrm{d}\\boldsymbol{\\alpha} = \\frac{1}{\\mathbb{Z}}\\int \\mathcal{L}(\\boldsymbol\\theta,\\boldsymbol{\\alpha})\\pi(\\boldsymbol\\theta,\\boldsymbol{\\alpha})\\,\\mathrm{d}\\boldsymbol{\\alpha}.\n\\end{equation}\nHere $\\mathbb{Z} \\equiv \\int \\mathcal{L}(\\boldsymbol\\theta,\\boldsymbol{\\alpha})\\pi(\\boldsymbol\\theta,\\boldsymbol{\\alpha})\\,\\mathrm{d}\\boldsymbol{\\alpha}\\,\\mathrm{d}\\boldsymbol{\\theta}$ is a normalisation factor referred to as the model evidence; taking ratios of evidences of different models is the most common method of Bayesian model comparison. In contrast, frequentist model testing typically involves building up the distribution of the likelihood or other test statistic by simulation, in order to determine the precise probability of obtaining the observed (or worse-fitting) data if the model is assumed to be correct.\n\nThe choice of Bayesian posterior or profile likelihood has strong implications for the required design of the algorithm with which to sample the model parameters: efficiently obtaining converged estimates of profile likelihood and posterior distributions requires drastically different sampling distributions. In neither case is random sampling at all sufficient nor correct, whether for accurate estimation of statistical properties nor for making statements about what is `typical' or `normal' within the parameter space of a given theory. Efficient profile likelihood evaluation requires fast location of and convergence towards the maximum likelihood, whereas efficient posterior and evidence evaluation requires samples obtained with a density approximately proportional to the value of the posterior itself (as indeed is the case in most other numerical integration problems).\n\nSampling and statistical considerations in \\textsf{GAMBIT}\\xspace are handled mostly by the \\textsf{ScannerBit}\\xspace module \\cite{ScannerBit}. It contains all the tools necessary to transform the likelihood function provided by the dependency resolver into converged profile likelihoods and Bayesian posteriors. It also facilitates Bayesian model comparison by calculating evidences (see e.g.\\ Refs.\\ \\cite{HP,Axions} for recent examples), and frequentist model testing by providing information that can be used to perform statistical simulations (see Ref.\\ \\cite{EWMSSM} for a detailed example).\n\nFor Bayesian analyses, \\textsf{ScannerBit}\\xspace provides a series of different prior transformations, allowing the user to choose what assumptions to make about the probabilities of different model parameter at the beginning of a run, and to sample accordingly.\n\n\\textsf{ScannerBit}\\xspace contains interfaces to various built-in and external implementations of a number of leading sampling algorithms. These include algorithms optimised for profile likelihood evaluation, and algorithms optimised for posterior and evidence calculations. Amongst these are \\twalk \\cite{ScannerBit}, a built-in ensemble Markov Chain Monte Carlo (MCMC) well suited to posterior evaluation, \\great \\cite{great}, a regular MCMC, \\diver \\cite{ScannerBit}, a differential evolution optimiser able to efficiently map profile likelihoods, and \\multinest \\cite{Feroz:2007kg,Feroz:2008xx} and \\textsf{polychord} \\cite{Handley:2015}, nested samplers well suited to evidence and posterior evaluation. Detailed performance comparisons between the different samplers can be found in Ref.\\ \\cite{ScannerBit}.\n\nFor consistency and the convenience of module function writers, \\textsf{GAMBIT}\\xspace also provides a series of relatively simple pre-profiled and pre-marginalised likelihood functions \\cite{gambit}. These functions provide likelihoods where the influence of one or more Gaussian or log-normally distributed nuisance parameters is profiled or integrated out without the assistance of explicit sampling by \\textsf{ScannerBit}\\xspace.\n\n\\subsection{Physics modules}\n\\label{sec:modules}\nThe physics content of \\textsf{GAMBIT}\\xspace currently resides in seven modules, which contain the module and backend functions relevant for all necessary theoretical calculations, simulations of particle astrophysics experiments and likelihood calculations. Future \\textsf{GAMBIT}\\xspace updates will both refine the code in each module, and add new modules for new branches of physics (such as the forthcoming \\textsf{CosmoBit} module).\n\n\\subsubsection{SpecBit}\n\nBSM physics theories necessarily introduce new particles. The first step in evaluating the likelihood of any parameter combination in a new theory is typically to calculate the masses and decay branching fractions of the new particles. These calculations get very complicated once one moves beyond tree level, as loop corrections can involve any number of new states in the theory, and loop corrections that shift the masses and decays of the existing SM particles must also be taken into account.\n\nParticle mass and coupling calculations are handled in the \\textsf{SpecBit}\\xspace module, which includes module functions for obtaining the pole masses and mixings of all new physical states in a model, scheme-dependent quantities such as those defined in the \\DR and $\\overline{MS}$\\xspace schemes, and SM masses and couplings (e.g. couplings of the SM-like Higgs). Generally, this information is obtained by running an appropriate spectrum generator but, in the simple case that the pole masses of a model are specified as input parameters, \\textsf{SpecBit}\\xspace simply formats the information to match that expected from a spectrum generator. In any case, it is important to realise that a spectrum cannot be stored simply as a set of numbers, since different experimental likelihoods may require predictions of running particle properties at different physical scales. Thus, \\textsf{SpecBit}\\xspace facilitates the passing of a spectrum object to module functions that contains knowledge of the renormalisation group equations of a model, allowing module functions in other parts of the \\textsf{GAMBIT}\\xspace code to locally run the \\DR or $\\overline{MS}$\\xspace parameters to the appropriate scale. Although \\textsf{SpecBit}\\xspace can be extended to include any model, its development has tended to proceed through updates that add functionality for the specific models targeted in \\textsf{GAMBIT}\\xspace physics papers. To date, this includes functions for GUT-\\cite{CMSSM} and weak-scale \\cite{MSSM,EWMSSM} parameterisations of the MSSM, singlet DM models with either a scalar, fermion or vector DM candidate \\cite{SSDM,SSDM2,HP}, minimal electroweak triplet and quintuplet DM \\cite{Piteration,McKay2}, and a low-energy object that holds SM-like particle information. A range of backends is used to supply the \\textsf{SpecBit}\\xspace calculations, including \\SPheno \\cite{Porod:2003um,Porod:2011nf} and \\FlexibleSUSY \\cite{Athron:2014yba} for BSM mass spectrum calculations. The latter is typically used for all spectrum generation requirements outside of the MSSM, including for the scalar singlet model examples described in this review. The Higgs and $W$ masses can also be calculated via the \\FeynHiggs \\cite{Heinemeyer:1998yj,Heinemeyer:1998np,Degrassi:2002fi,Frank:2006yh,Hahn:2013ria,Bahl:2016brp} backend.\n\n\\subsubsection{DecayBit}\n\nParticle decay calculations are handled by the \\textsf{DecayBit}\\xspace module, after accepting the masses and couplings of particles from \\textsf{SpecBit}\\xspace. These are used to calculate decay widths and branching fractions for each particle, which are stored in a single decay table entry for each particle. The collection of entries is then gathered into a full \\gambit decay table, which is passed on to other \\gambit modules.\n\n\\textsf{DecayBit}\\xspace includes known SM particle decays, modifications of SM particle decays through new physics effects, and the decays of BSM particles. For the SM, \\textsf{DecayBit}\\xspace contains the Particle Data Group results for the total widths for the $W$, $Z$, $t$, $b$, $\\tau$ and $\\mu$ (plus antiparticles), and for the most common mesons $\\pi^0$, $\\pi^\\pm$, $\\eta$, $\\rho^0, \\rho^\\pm$ and $\\omega$~\\cite{PDB}. In addition, partial widths to all distinct final states are provided for $W$, $Z$, $t$, $b$, $\\tau$, $\\mu$, $\\pi^0$ and $\\pi^\\pm$. These ``pure SM'' decays are used in \\textsf{GAMBIT}\\xspace whenever an SM decay acquires no BSM contribution in a model, or when the only effect of the BSM physics is to introduce a new decay channel, in which case the pure SM decays can be appended to the new list of decay channels. For the pure SM Higgs boson, the user can decide whether to calculate the partial and total decay widths at the predicted value of the Higgs mass with \\textsf{FeynHiggs}\\xspace, or to use pre-computed tables provided in \\textsf{DecayBit}\\xspace, sourced from Ref.\\ \\cite{YellowBook13}.\n\nBSM decays are handled on a model-by-model basis. For Higgs portal DM models, \\textsf{DecayBit}\\xspace contains analytic expressions for the partial width for a Higgs decay to two DM particles, and this decay is added to the list of SM Higgs partial widths, before rescaling the decay branching fractions and the total width. For MSSM variants, \\textsf{DecayBit}\\xspace calculates both the decays of all sparticles and additional Higgs bosons, and the SUSY corrections to the decays of the SM-like Higgs boson and the top quark. Higgs decay results may be sourced from either \\HDECAY via \\SUSYHIT, or \\FeynHiggs, whilst top quark decays are only available via \\FeynHiggs. Sparticle decays are obtained from \\SDECAY via \\SUSYHIT, but we note that a patch to the code is required to allow \\gambit to call functions from a shared library, and to solve problems with negative decay widths for some models due to large and negative 1-loop QCD corrections. Full details are given in~\\cite{SDPBit}; the patch is applied automatically when \\SUSYHIT is retrieved and built from within \\textsf{GAMBIT}\\xspace.\n\nA recent update of \\textsf{DecayBit}\\xspace has seen the addition of routines for observables relating to right-handed neutrino studies. This includes the invisible decay width of the $Z$ boson $\\Gamma_{\\rm{inv}}$, and the leptonic $W$ boson decay widths $\\Gamma_{W\\to e \\bar\\nu_e}$, $\\Gamma_{W\\to \\mu \\bar\\nu_\\mu}$ and $\\Gamma_{W\\to \\tau \\bar\\nu_\\tau}$. Measurements and uncertainties are taken from Ref.~\\cite{PDG17}, whilst theoretical results are taken from Refs.~\\cite{Drewes:2015iva,Antusch:2014woa,Antusch:2015mia,Ferroglia:2012ir,Antusch:2015mia,Fernandez-Martinez:2015hxa,Abada:2013aba,Dubovyk:2018rlg,Antusch:2014woa}.\n\n\\subsubsection{PrecisionBit}\n\nSome of the most severe constraints on BSM physics scenarios come from precision measurements of the electroweak sector, and other SM quantities. In \\textsf{GAMBIT}\\xspace, these are handled by the \\textsf{PrecisionBit}\\xspace module, which provides nuisance likelihoods for SM quantities such as the top quark mass and strong coupling constant, which have been measured with high precision. A related function is the calculation of the BSM corrections to SM observables such as the mass of the $W$ boson and the weak mixing angle, and the provision of likelihood functions that compare these predictions with experimental data.\n\nA schematic representation of \\textsf{PrecisionBit}\\xspace is shown in Fig.\\ \\ref{fig:precisionbit}, providing our first interesting example of the interaction between \\textsf{GAMBIT}\\xspace modules. Standard Model parameters that do not require BSM correction calculations are provided directly by the \\textsf{GAMBIT}\\xspace SM model, and are used in the calculation of SM nuisance likelihoods. The BSM parameters, meanwhile, are first used by \\textsf{SpecBit}\\xspace in the calculation of particle masses, couplings and precision Higgs properties. \\textsf{PrecisionBit}\\xspace then updates the results to form a precision-updated spectrum (including a dedicated calculation of the $W$ mass) which is used for calculating Higgs and $W$ mass likelihoods, in addition to a suite of electroweak precision observables (EWPO).\n\nLikelihoods exist for the Fermi coupling constant ($G_F$), the fine-structure constant ($\\alpha_{\\mathrm{em}}$), the $\\overline{MS}$\\xspace light quark ($u$,$d$,$s$) masses at $\\mu=2$~GeV, the charm ($m_c(m_c)$) and bottom ($m_b(m_b)$) masses, and the $W$, $Z$ and Higgs boson masses. There are also calculations and likelihoods for other precision observables such as the anomalous magnetic moment of the muon $a_\\mu=\\frac12(g-2)_\\mu$, the effective leptonic weak mixing angle sin$^2\\theta_{W,eff}$, and the departure from 1 of the ratio of the Fermi constants implied by the neutral and weak currents $\\Delta\\rho$. Note that, for the full suite of observables, calculations are currently only included for the MSSM; calculations for other models will be added as the corresponding models are implemented in \\textsf{GAMBIT}\\xspace.\n\n\\begin{figure}[t]\n\\centering\n\\includegraphics[width=0.59\\textwidth]{figures\/PrecisionBit}\n\\caption{Schematic representation of the structure of \\textsf{PrecisionBit}\\xspace. From \\cite{SDPBit}.}\n\\label{fig:precisionbit}\n\\end{figure}\n\nLike \\textsf{DecayBit}\\xspace, \\textsf{PrecisionBit}\\xspace has also recently been updated to include observables for right-handed neutrino studies. These consist of right-handed neutrino contributions to $m_W$ and sin$^2\\theta_{W,eff}$.\n\n\\subsubsection{DarkBit}\nBSM physics models that include particle DM candidates can potentially give rise to observable consequences in a wide range of astrophysical DM experiments.\n\nIn gamma-ray indirect detection, the \\textsf{DarkBit}\\xspace module contains a dedicated signal yield calculator, along with an interface to \\gamlike, a likelihood calculator for current and future gamma-ray experiments. This combination can cope with signatures that result from an arbitrary mixture of final states, which significantly extends previous tools.\n\nFurther indirect detection constraints come from an interface to the \\nulike neutrino telescope likelihood package~\\cite{IC79_SUSY}.\n\nDirect DM search experiments are handled by the dedicated \\ddcalc package, which can be extended to include the effects of generic interactions between Weakly Interacting Massive Particles (WIMPs) and nucleons, as parameterised through effective operators. This includes both spin-dependent and spin-independent scattering. The package models a wide range of direct search experiments including Xenon100, SuperCDMS, SIMPLE, LUX, PandaX, PICO-60 and PICO-2L.\n\nFinally, the relic density of dark matter can be computed via interfaces to \\textsf{DarkSUSY}\\xspace and \\textsf{micrOMEGAs}\\xspace~\\cite{darksusy,micromegas_nu}, and used to constrain models by computing a likelihood based on the value observed by \\textit{Planck} \\cite{Planck18cosmo}.\n\nThe basic structure of \\textsf{DarkBit}\\xspace applicable to WIMP theories is sketched in Fig.\\ \\ref{fig:flow}, providing a good example of \\textsf{GAMBIT}\\xspace's modular design principle. None of the likelihoods requires knowledge of the BSM physics parameters, instead only requiring knowledge of derived quantities that can be shared between likelihood calculations. The first step in \\textsf{DarkBit}\\xspace is to create a Process Catalogue containing information on particle annihilation processes, using the particle masses and couplings provided by \\textsf{SpecBit}\\xspace. For indirect detection calculations, this is used to create the gamma ray or neutrino spectrum of the annihilation products, via a weighted sum of indiviual contributions. For long decay chains, a native cascade decay Monte Carlo generator is used. This final annihilation spectrum is then passed to the likelihood calculators for gamma ray and neutrino telescope experiments. The Process Catalogue is also used to provide the effective annihilation rate for relic density calculations, which is then passed to a Boltzmann solver, followed by the relic density likelihood calculator. For direct dectection signatures, the model parameters are used to set the WIMP-nucleon couplings, which are then used in the calculation of the direct detection likelihood via the \\ddcalc package.\n\nA recent update of \\textsf{DarkBit}\\xspace has added various module functions required for the calculation of axion observables and likelihoods. The included observables are detailed in Section~\\ref{sec:axionL}.\n\n\\begin{figure*}[t]\n \\centering\n \\includegraphics[width=0.60\\linewidth]{figures\/DarkBit_flow}\n \\caption{Schematic overview of the \\textsf{DarkBit}\\xspace module. The two-letter insets indicate what backend codes can\n be used: \\textsf{DarkSUSY}\\xspace (DS), \\textsf{micrOMEGAs}\\xspace (MO), \\gamLike (GL), \\nulike (NL) and\n \\ddcalc (DC). From \\cite{DarkBit}.\n }\n \\label{fig:flow}\n\\end{figure*}\n\nAn important additional function of \\textsf{DarkBit}\\xspace is to constrain nuisance parameters for various astrophysical unknowns that strongly affect direct and indirect searches for DM. \\textsf{DarkBit}\\xspace contains likelihoods for the parameters of the local DM spatial and velocity distributions, plus the nuclear matrix elements that enter direct search WIMP-nucleon scattering calculations.\n\n\\subsubsection{FlavBit}\n\nA very powerful indirect probe of BSM physics comes from the measurement of flavour physics processes, as theoretical predictions for these observables would be shifted by loop corrections from new particles. The excellent precision of flavour phyics measurements allows them to be sensitive to much higher energy scales than direct searches for new particles. Indeed, recent measurements from the LHCb experiment~\\cite{Aaij:2016flj,Aaij:2015esa,Aaij:2015yra,Aaij:2014ora,Aaij:2013qta,Aaij:2019wad,Aaij:2014pli} and from $B$ factories~\\cite{Lees:2012xj,Lees:2013uzd,Huschle:2015rga,Abdesselam:2016cgx,Abdesselam:2016llu,Aubert:2006vb,Lees:2015ymt,Wei:2009zv,Wehle:2016yoi} show tensions with the SM that are generating a considerable amount of theoretical interest.\n\n\\textsf{FlavBit}\\xspace implements flavour physics constraints from rare decay observables using the effective Hamiltonian approach, in which the cross-sections for transitions from initial states $i$ to final states $f$ are proportional to the squared matrix elements $|\\langle f |{\\cal H}_{\\rm eff}|i\\rangle|^2$. For example, an effective Hamiltonian for $b \\rightarrow s$ transitions given by\n\\begin{equation}\n\\mathcal{H}_{\\rm eff} = -\\frac{4G_{F}}{\\sqrt{2}} V_{tb} V_{ts}^{*} \\sum_{i=1}^{10} \\Bigl(C_{i}(\\mu) \\mathcal{O}_i(\\mu)+C'_{i}(\\mu) \\mathcal{O}'_i(\\mu)\\Bigr)\\;.\n\\end{equation}\nThe local operators $\\mathcal{O}_i$ represent long-distance interactions. The Wilson coefficients $C_i$ can be calculated using perturbative methods, by requiring matching between the high-scale theory and the low-energy effective theory, at some scale $\\mu_W$ which is of the order of $m_W$. The Wilson coefficients can then be evolved to the characteristic scale for $B$ physics calculations ($\\mu_b$, of the order of the $b$ quark mass) using the renormalisation group equations of the \\emph{effective} field theory. A similar approach can be taken to $b\\rightarrow d$ transitions, using a different basis of low-energy operators. The original list of observables in \\textsf{FlavBit}\\xspace was divided into four categories:\n\n\\begin{itemize}\n\\item {\\bf Tree-level leptonic and semi-leptonic decays: }includes decays of $B$ and $D$ mesons to leptons, including $B^\\pm \\to \\tau \\nu_\\tau$, $B \\to D^{(*)} \\tau \\nu_\\tau$ and $B \\to D^{(*)} \\ell \\nu_\\ell$.\n\\item {\\bf Electroweak penguin transitions: }includes measurements of rare decays of the form $B \\to M \\ell^+\\ell^-$ (where $M$ is a meson with a smaller mass than the parent meson), such as angular observables of the decay $B^0 \\to K^{*0} \\mu^+\\mu^-$.\n\\item {\\bf Rare purely leptonic decays: }includes $B$ decays with only leptons in the final state, such as $B^0_{(s)} \\to \\mu^+ \\mu^-$.\n\\item {\\bf Other observables: }includes $b\\to s$ transitions in the radiative decays $B \\to X_s \\gamma$, the mass difference ($\\Delta M_s$ between the heavy $B_H$ and light $B_L$ eigenstates of the $B^0_s$ system, and kaon and pion decays (e.g. the leptonic decay ratio ${\\cal B}(K^\\pm\\to \\mu \\nu_\\mu)\/{\\cal B}(\\pi^\\pm\\to \\mu \\nu_\\mu)$).\n\\end{itemize}\n\nTheoretical calculations for these processes are handled via an interface to \\superiso~\\cite{Mahmoudi:2007vz,Mahmoudi:2008tp,Mahmoudi:2009zz}. Experimental results used in the calculation of likelihoods come from a variety of sources, including the PDG, the BaBar and Belle experiments, the HFAG collaboration and the LHCb experiment. Full details are given in~\\cite{FlavBit}.\n\nMore recently, \\textsf{FlavBit}\\xspace has been updated with observables relevant to right-handed neutrinos. These include:\n\\begin{itemize}\n\\item {\\bf Lepton-flavour violating (LFV) muon and tau decay searches} performed by the MEG, BaBar, Belle, ATLAS, LHCb and SINDRUM collaborations~\\cite{TheMEG:2016wtm,Aubert:2009ag,Hayasaka:2007vc,Lees:2010ez,Hayasaka:2010np,Aad:2016wce,Aaij:2014azz,Bellgardt:1987du}. LFV processes can also result in a neutrinoless $\\mu - e$ conversion inside a nucleus, and these are included in the form of three results using Ti, Pb and Au nuclei obtained by the SINDRUM II experiment~\\cite{Kaulard:1998rb,Honecker:1996zf,Bertl:2006up}.\n\\item {\\bf Tests of lepton universality violation} in the semileptonic decays of $B$ mesons $B^{0\/\\pm} \\to X^{0\/\\pm} l^+ l^-$, as performed by LHCb \\cite{Aaij:2014ora,Aaij:2017vbb}.\n\\end{itemize}\n\nA forthcoming major update to the \\textsf{FlavBit}\\xspace module will add an interface for \\superisofour, with added support for theory uncertainty covariance matrices. The experimental likelihoods will also receive an update, via a new interface to the \\heplike package \\cite{HEPLike}.\n\n\\begin{figure}[tp]\n\\centering\n\\includegraphics[width=0.60\\linewidth]{figures\/colliderbit-flow}\n\\caption{Schematic diagram of the \\textsf{ColliderBit}\\xspace processing chain for LHC likelihoods. From \\cite{ColliderBit}.}\n\\label{fig:lhcchain}\n\\end{figure}\n\n\\subsubsection{ColliderBit}\n\\label{sec:colliderbit}\nA leading source of constraints on BSM physics models comes from high-energy collider searches for new particles, plus the relatively recent measurements of the Higgs boson mass and decay branching fractions. The \\textsf{ColliderBit}\\xspace module includes the most comprehensive list of recent LHC particle searches of any public package, alongside a new interpolation of LEP results for supersymmetric particle searches. Higgs signal strength and mass measurements (including limits on possible signatures arising from new Higgs bosons) are handled via an interface to the \\textsf{HiggsSignals}\\xspace~\\cite{HiggsSignals} and \\textsf{HiggsBounds}\\xspace~\\cite{Bechtle:2008jh,Bechtle:2011sb} packages, which includes data from LEP, the Tevatron and the LHC.\n\nLHC constraints are particularly difficult to model rigorously for general models. Searches for new particles are often optimised on, and interpreted in terms of, so-called ``simplified models'', which feature only a few options from the much broader phenomenology of the parent model. For example, searches for supersymmetric particles might assume that only a particular pair of sparticles is ever produced, with decays fixed to a particular final state. The resulting exclusion limit will never apply directly to a more general model, although one can obtain approximate limits by scaling individual simplified model limits by the known cross-sections and branching ratios for each parameter point~~\\cite{Kraml:2013mwa,Papucci:2014rja}.\n\nIn \\textsf{ColliderBit}\\xspace, we provide more rigorous limits by performing an actual reproduction of the ATLAS and CMS limit-setting procedures, as shown in Fig.\\ \\ref{fig:lhcchain}. This includes a cross-section calculation for new particle production processes, followed by Monte Carlo simulation of LHC events for each parameter point using a custom parallelised version of the \\textsf{Pythia\\,8}\\xspace generator~\\cite{Sjostrand:2006za,Sjostrand:2014zea}. The results can either be fed at the truth level into code that reproduces the kinematic selections of a wide range of LHC analyses, or passed through a custom detector simulation based on four-vector smearing before analysis. Cross-sections are currently taken at leading order (plus leading log) from \\textsf{Pythia\\,8}\\xspace, but a forthcoming update will allow user-specified cross-sections. The final step of the process is to calculate a combined likelihood by either taking the signal region in a given final state for each experiment that is expected to have the highest sensitivity to the model in question, or by using a covariance matrix for analyses in cases where this is published by the relevant experimental collaboration. The list of \\textsf{ColliderBit}\\xspace analyses is continually updated, and currently includes a broad selection of searches for supersymmetric particles, plus monojet searches for DM particles.\n\n\\subsubsection{NeutrinoBit}\n\\label{sec:neutrinobit}\nThe \\textsf{NeutrinoBit}\\xspace module contains a variety of module functions for calculating observables and likelihoods in the neutrino sector, both for SM(-like) neutrinos and for right-handed neutrinos (RHNs). RHNs could cause observable consequences in a number of experiments, although it should be noted that the recent \\textsf{GAMBIT}\\xspace study focussed on models that are capable of explaining the light neutrino oscillation data, which excludes most sterile neutrino dark matter models. This is because long-lived RHNs would require very small couplings with SM matter, in which case their contribution to light neutrino mass generation is negligible.\n\n\\textsf{NeutrinoBit}\\xspace currently contains likelihoods dealing with the following classes of experimental data:\n\n\\begin{itemize}\n\\item \\textbf{Active neutrino mixing: }\\textsf{NeutrinoBit}\\xspace includes likelihoods for the 3-flavour SM-like active neutrino mixing observables $\\theta_{12}$, $\\theta_{13}$, $\\theta_{23}$ (mixing angles), $\\delta_{\\mathrm{CP}}$ (CP-phase) and $\\Delta m^2_{21}$ and $\\Delta m^2_{3\\ell}$ (mass splittings) with $\\ell = 1$ for normal mass ordering and $\\ell = 2$ for inverted mass ordering. The likelihoods use the one-dimensional $\\Delta\\chi^2$ tables provided by the NuFIT collaboration~\\cite{Esteban:2016qun,NuFit}. These in turn include results from the solar neutrino experiments Homestake (chlorine)~\\cite{Cleveland:1998nv}, Gallex\/GNO~\\cite{Kaether:2010ag}, SAGE~\\cite{Abdurashitov:2009tn}, SNO~\\cite{Aharmim:2011vm}, the four phases of Super-Kamiokande~\\cite{Hosaka:2005um,Cravens:2008aa,Abe:2010hy} and two phases of Borexino~\\cite{Bellini:2011rx,Bellini:2008mr,Bellini:2014uqa}. They also include results from the atmospheric experiments IceCube\/DeepCore~\\cite{Aartsen:2014yll}, the reactor experiments KamLAND~\\cite{Gando:2013nba}, Double-Chooz~\\cite{An:2016srz}, Daya-Bay~\\cite{An:2016ses} and Reno~\\cite{reno}, the accelerator experiments MINOS~\\cite{Adamson:2013whj,Adamson:2013ue}, T2K~\\cite{t2k} and NO$\\nu$A~\\cite{nova}, and the cosmic microwave background results from Planck~\\cite{Ade:2015xua}.\n\\item \\textbf{Lepton universality: }\\textsf{NeutrinoBit}\\xspace contains likelihoods for lepton universality violation in fully leptonic decays of charged mesons, $X^+ \\to l^+ \\nu$.\n\\item \\textbf{CKM unitarity: }The determination of the CKM matrix elements usually relies on the assumption that the active-sterile neutrino mixing matrix is zero. The presence of non-trivial mixing thus modifies the CKM matrix elements, and the experimentally-observed values can be used to simultaneously constrain the true CKM element values, and the active-sterile mixing matrix $\\Theta$. \\textsf{NeutrinoBit}\\xspace constructs a likelihood based on the deviations of the true values of $(V_{CKM})_{us}$ and $(V_{CKM})_{ud}$ from their experimentally-measured values.\n\\item \\textbf{Neutrinoless double-beta decay: }In a double-beta decay process, two neutrons decay into two protons, with the emission of two electrons and two anti-neutrinos. Majorana neutrinos would give rise to lepton number violation, resulting in neutrinoless double-beta decay ($0\\nu\\beta\\beta$). In addition, the exchange of RHNs can modify the effective neutrino mass $m_{\\beta\\beta}$, which is constrained by half-life measurements of $0\\nu\\beta\\beta$ decay. The best upper limits currently come from the GERDA experiment (Germanium)~\\cite{Agostini:2017iyd} with $m_{\\beta\\beta}<0.15-0.33\\;\\text{eV}$ (90\\% CL), and KamLAND-Zen (Xenon)~\\cite{KamLAND-Zen:2016pfg}, $m_{\\beta\\beta}<0.061-0.165\\;\\text{eV}$ (90\\% CL). \\textsf{NeutrinoBit}\\xspace uses these values to define one-sided Gaussian likelihoods, with theoretical calculations for RHN models taken from Refs.~\\cite{Drewes:2016lqo,Faessler:2014kka}\n\\item \\textbf{Big Bang Nucleosynthesis: }RHNs can affect the abundances of the primordial elements if they decay shortly before BBN, as the typical energy of the decay products is significantly higher than the plasma energy at that time. This can lead to the dissociation of formed nuclei, or the creation of deviations from thermal equilibrium. The requirement that RHNs decay before BBN implies an upper limit on their lifetime which, in turn, results in a constraint on the total mixing with the active neutrinos. \\textsf{NeutrinoBit}\\xspace currently includes a basic BBN likelihood that uses decay expressions from Refs.~\\cite{Gorbunov:2007ak,Canetti:2012kh}, and requires the lifetime of each RHN to be less than 0.1s~\\cite{Ruchayskiy:2012si}. A more comprehensive update will be released in future, associated with the new \\textsf{CosmoBit} module.\n\\item \\textbf{Direct RHN searches: }Direct searches for RHNs can be performed by looking for peaks in the lepton energy spectrum of a meson decay, looking for evidence of production in beam dump experiments, and by studying the decay of vector bosons or mesons in $e^+e^-$ or $pp$ colliders. \\textsf{NeutrinoBit}\\xspace contains likelihoods for RHN searches at the PIENU~\\cite{PIENU:2011aa}, PS-191~\\cite{Bernardi:1987ek}, CHARM~\\cite{Bergsma:1985is}, E949~\\cite{Shaykhiev:2011zz,Artamonov:2014urb}, NuTeV~\\cite{Vaitaitis:1999wq}, DELPHI~\\cite{Abreu:1997uq}, ATLAS~\\cite{Aad:2015xaa} and CMS~\\cite{Sirunyan:2018mtv} experiments.\n\\end{itemize}\n\n\\section{Applications to new physics}\n\\label{sec:physics}\n\\subsection{Supersymmetry}\n\\label{sec:SUSY}\n\nSupersymmetry (SUSY) has long been one of the leading candidates for BSM physics, owing to its potential for simultaneously answering several of the questions left open by the SM. In particular, the hierarchy problem and the dark matter ``WIMP miracle'' suggest the possible existence of SUSY states around the weak scale.\n\nMost phenomenological explorations of SUSY take the MSSM as their starting point. On top of its minimal supersymmetrisation of the SM, the MSSM effectively parameterises our ignorance about the high-scale mechanism responsible for breaking SUSY. This is done by including in the Lagrangian all gauge-invariant and renormalisable terms that break SUSY ``softly'', that is, without re-introducing the quadratic Higgs mass divergences that gave rise to the hierarchy problem. In this way the MSSM provides a unified framework for exploring a wide range of possible manifestations of SUSY, but at the price of a vast parameter space: if no further assumptions are made the soft SUSY-breaking terms introduce more than one hundred free parameters.\n\nMany different assumptions have been employed in the literature to reduce this parametric freedom and improve predictability. The resulting models broadly fit in two categories.\n\nThe first category consists of high-scale models that take inspiration from the fact that SUSY can provide gauge coupling unification at some high Grand Unified Theory (GUT) scale, typically around $10^{16}$\\,GeV. In these models a small number of unified mass and coupling parameters are defined at the GUT scale and then run down to the electroweak scale where phenomenological predictions are calculated. Thus, the assumption of high-scale unification constrains the model to a low-dimensional subspace of the full MSSM space, effectively imposing a set of characteristic correlations among the many MSSM parameters at the weak scale.\n\nProbably the most studied SUSY model in this category is the Constrained MSSM (\\textsf{CMSSM}\\xspace)~\\cite{Nilles:1983ge}. Here the parameter space is reduced to only four continuous parameters and a sign choice: the unified soft-breaking scalar mass, $m_0$; the unified soft-breaking gaugino mass, $m_{1\/2}$; the unified trilinear coupling, $A_0$; the ratio of the vacuum expectation values for the two Higgs doublets, $\\tan\\beta\\equiv v_\\mathrm{u}\/v_\\mathrm{d}$; and the sign of the supersymmetric Higgsino mass parameter $\\mu$. The \\textsf{CMSSM}\\xspace has been studied in global fits for over a decade~\\cite{Han:2016gvr, Bechtle:2014yna, arXiv:1405.4289, arXiv:1402.5419, MastercodeCMSSM, arXiv:1312.5233, arXiv:1310.3045, arXiv:1309.6958, arXiv:1307.3383, arXiv:1304.5526, arXiv:1212.2886, Strege13, Gladyshev:2012xq, Kowalska:2012gs, Mastercode12b, arXiv:1207.1839, arXiv:1207.4846, Roszkowski12, SuperbayesHiggs, Fittino12, Mastercode12, arXiv:1111.6098, Fittino, Trotta08, Ruiz06, Allanach06, Fittino06, Baltz04, SFitter}, most recently in the \\textsf{GAMBIT}\\xspace analysis in~\\cite{CMSSM}.\n\n\nTwo much-studied generalisations of the \\textsf{CMSSM}\\xspace are the Non-Universal Higgs Mass models 1 and 2 (\\textsf{NUHM1}\\xspace and \\textsf{NUHM2}\\xspace)~\\cite{Matalliotakis:1994ft,Olechowski:1994gm,Berezinsky:1995cj,Drees:1996pk,Nath:1997qm}. These models loosen the tight link in the \\textsf{CMSSM}\\xspace between the Higgs sector and the sfermions by separating out the soft-breaking mass parameters of the Higgs sector from the common scalar mass parameter $m_0$. This is achieved by introducing either one (\\textsf{NUHM1}\\xspace) or two (\\textsf{NUHM2}\\xspace) additional parameters at the GUT scale. In recent years the \\textsf{NUHM1}\\xspace and \\textsf{NUHM2}\\xspace have been studied in several global fit analyses~\\cite{MastercodeCMSSM, arXiv:1312.5233, Strege13, Fittino12, Mastercode15, Buchmueller:2014yva, arXiv:1405.4289}. \\textsf{GAMBIT}\\xspace global fits of the \\textsf{NUHM1}\\xspace and \\textsf{NUHM2}\\xspace were performed along with the fit of the \\textsf{CMSSM}\\xspace in~\\cite{CMSSM}. In Section~\\ref{sec:GUT-SUSY} we summarise the \\textsf{GAMBIT}\\xspace results for these GUT-scale SUSY models.\n\nThe second category of MSSM sub-models are the weak-scale models. Here the focus is on exploring a broad range of weak-scale phenomenological scenarios in an economical manner, by varying only the MSSM parameters that most directly impact the observables under study. With all MSSM parameters defined near the weak scale, these models are mostly agnostic to questions concerning physics at very high scales, such as grand unification. The models are often labeled as \\textsf{MSSM}$n$ (or as \\textsf{pMSSM}$n$ for the \\textit{phenomenological} \\textsf{MSSM}$n$), with $n$ specifying the number of weak-scale MSSM parameters that are treated as free parameters.\n\nVarious such weak-scale models have been subjected to global fit analyses in the past few years~\\cite{arXiv:1608.02489, arXiv:1507.07008, Mastercode15, arXiv:1506.02499, arXiv:1504.03260, Mastercode17}. The \\textsf{GAMBIT}\\xspace analyses in this category are~\\cite{MSSM}, which looks at a seven-dimensional MSSM parameterisation (\\textsf{MSSM7}\\xspace), and~\\cite{EWMSSM}, in which the fast LHC simulation capabilities of \\textsf{ColliderBit}\\xspace are used for a collider-focused fit of the four-dimensional MSSM ``electroweakino'' (chargino and neutralino) sector (\\textsf{EWMSSM}\\xspace). We summarise the \\textsf{GAMBIT}\\xspace results for the \\textsf{MSSM7}\\xspace in Section~\\ref{sec:MSSM7} and for the \\textsf{EWMSSM}\\xspace in Section~\\ref{sec:EWMSSM}.\n\nThe phenomenological richness of the MSSM means that a wide range experimental results are relevant for constraining the parameter space. The mass and signal strength measurements for the $125$\\,GeV Higgs boson and the measurement of the relic density of dark matter are of particular importance. We note that the impact of the relic density measurement depends strongly on whether the SUSY model is assumed to account for the full relic density, or, more conservatively, some arbitrary fraction of it. The \\textsf{GAMBIT}\\xspace studies reviewed here all take the latter approach.\n\nMeasurements of electroweak precision observables such as $m_W$ and the muon $g-2$, and of flavour observables like $BR(B \\rightarrow X_s \\gamma)$ and $BR(B_{(s)} \\rightarrow \\mu^+ \\mu^-)$, introduce further important requirements on the SUSY parameter space. Finally, the null results from direct and indirect dark matter searches, and from collider searches for sparticles and additional Higgs bosons, essentially rule out some parts of SUSY parameter space. Though, to determine the exact implications of such null-result collider searches is far from trivial, as will be illustrated by the \\textsf{EWMSSM}\\xspace results discussed in Section~\\ref{sec:EWMSSM}.\n\n\n\\subsubsection{Results for the \\textsf{CMSSM}\\xspace, \\textsf{NUHM1}\\xspace and \\textsf{NUHM2}\\xspace}\n\\label{sec:GUT-SUSY}\n\nThe \\textsf{GAMBIT}\\xspace global fits of the \\textsf{CMSSM}\\xspace, \\textsf{NUHM1}\\xspace and \\textsf{NUHM2}\\xspace in~\\cite{CMSSM} are interpreted in terms of frequentist profile likelihood maps, identifying the best-fit point and the $1\\sigma$ and $2\\sigma$ preferred regions relative to this point. As the results for \\textsf{NUHM2}\\xspace are qualitatively similar to those for \\textsf{NUHM1}\\xspace, we here focus on the \\textsf{CMSSM}\\xspace and \\textsf{NUHM1}\\xspace results.\n\n\\begin{figure}\n \\centering\n \\includegraphics[width = 0.49\\textwidth]{figures\/susy\/CMSSM_2_3_like2D_post}\n \\includegraphics[width = 0.49\\textwidth]{figures\/susy\/CMSSM_2_3_ColourMechanism_post}\\\\\n \\includegraphics[width = 0.49\\textwidth]{figures\/susy\/NUHM1_2_3_like2D_post}\n \\includegraphics[width = 0.49\\textwidth]{figures\/susy\/NUHM1_2_3_ColourMechanism_post}\\\\\n \\includegraphics[height=3.1mm]{figures\/susy\/rdcolours4.pdf}\n \\caption{\n Profile likelihood in the $(m_0,m_{1\/2})$ plane in the \\textsf{CMSSM}\\xspace (\\textit{top left}) and the \\textsf{NUHM1}\\xspace (\\textit{bottom left}). The right-hand panels show the mechanisms that contribute to bringing the predicted DM relic density close to or below the observed value. The white contours show the $1\\sigma$ and $2\\sigma$ preferred regions relative to the best-fit point (white star). From~\\cite{CMSSM}.\n }\n \\label{fig:CMSSM_NUHM1_m0_m12}\n\\end{figure}\n\nThe profile likelihood maps for the $(m_0,m_{1\/2})$ planes of the \\textsf{CMSSM}\\xspace and \\textsf{NUHM1}\\xspace are shown in the left panels of Fig.\\ \\ref{fig:CMSSM_NUHM1_m0_m12}. The \\textsf{NUHM1}\\xspace plane is clearly less constrained compared to the \\textsf{CMSSM}\\xspace. The underlying reason is the additional parametric freedom in the Higgs sector of the \\textsf{NUHM1}\\xspace, where the MSSM Higgs parameters $m_{H_u}^2$ and $m_{H_d}^2$ are not unified with $m_0^2$ at the GUT scale, but are rather set by an independent parameter $m_H$ through the GUT-scale requirement $m_{H_u} = m_{H_d} \\equiv m_H$. (In the \\textsf{NUHM2}\\xspace, $m_{H_u}$ and $m_{H_d}$ are taken as independent parameters at the GUT scale.) We note that as $m_H$ is taken to be a real parameter, we have $m_{H_u}^2 = m_{H_d}^2 > 0$ at the GUT scale. The correct shape of the Higgs potential at the weak scale must therefore be generated through radiative corrections, as is the case for the \\textsf{CMSSM}\\xspace.\n\nThe right-hand panels in Fig.\\ \\ref{fig:CMSSM_NUHM1_m0_m12} help us understand the preferred parameter space in more detail. In these panels different sub-regions of the $2\\sigma$ region are coloured according to which mechanism(s) contribute to keeping the DM relic density close to or below the observed value. The following criteria are used to define the DM mechanism regions in~\\cite{CMSSM} and in the \\textsf{MSSM7}\\xspace study in~\\cite{MSSM}:\n\\begin{itemize}\n\\item stop co-annihilation: $m_{\\tilde{t}_1} \\leq 1.2\\,m_{\\tilde{\\chi}^0_1}$,\n\\item sbottom co-annihilation: $m_{\\tilde{b}_1} \\leq 1.2\\,m_{\\tilde{\\chi}^0_1}$,\n\\item stau co-annihilation: $m_{\\tilde{\\tau}_1} \\leq 1.2\\,m_{\\tilde{\\chi}^0_1}$,\n\\item chargino co-annihilation: $\\tilde{\\chi}^0_1$ $\\ge50\\%$ Higgsino,\\footnote{For brevity we refer to this mechanism simply as ``chargino co-annihilation'', though it also includes co-annihilations with the next-to-lightest neutralino. Further, for many points in this region the most important effect is simply enhanced $\\tilde{\\chi}^0_1$--$\\tilde{\\chi}^0_1$ annihilations, owing to the dominantly Higgsino $\\tilde{\\chi}^0_1$ composition.}\n\\item $A\/H$ funnel: $1.6\\,m_{\\tilde{\\chi}^0_1} \\leq \\textrm{($m_A$ or $m_H$)} \\leq 2.4\\,m_{\\tilde{\\chi}^0_1}$,\n\\item $h\/Z$ funnel: $1.6\\,m_{\\tilde{\\chi}^0_1} \\leq \\textrm{($m_Z$ or $m_h$)} \\leq 2.4\\,m_{\\tilde{\\chi}^0_1}$.\n\\end{itemize}\nThe coloured regions overlap for parameter points where more than one mechanism contributes.\n\nIn the \\textsf{CMSSM}\\xspace the overall highest-likelihood point is found in the stop co-annihilation region, at $m_0 \\lesssim 4.5$\\,TeV. This region is associated with large, negative values for the trilinear coupling, $A_0 \\lesssim -5$\\,TeV, and $\\tan\\beta \\lesssim 16$. Only two other DM mechanisms are active within the best-fit parameter space, namely the $A\/H$ funnel and chargino co-annihilation. Thus, in contrast with earlier \\textsf{CMSSM}\\xspace fits, these results show that the stau co-annihilation region has fallen out of the preferred parameter space. This is mainly driven by the likelihood contribution from the LHC Higgs measurements, which penalise the lower-$m_0$ region where the lightest stau get sufficiently close in mass to the lightest neutralino.\n\nAs discussed above, the link between $m_0$ and the Higgs sector is relaxed in the \\textsf{NUHM1}\\xspace. This opens up the parameter space at lower $m_0$, allowing the stau co-annihilation region back within the $2\\sigma$ preferred region, as seen in the lower right panel of Fig.\\ \\ref{fig:CMSSM_NUHM1_m0_m12}. A second consequence of $m_0$ being decoupled from the Higgs sector in the \\textsf{NUHM1}\\xspace is that the allowed chargino co-annihilation region is extended to arbitrarily small $m_0$ values, compared to in the \\textsf{CMSSM}\\xspace. We can understand this by investigating the \\textsf{CMSSM}\\xspace case: the chargino co-annihilation DM mechanism is important when the MSSM Higgsino mass parameter $\\mu$ is smaller than the bino mass parameter $M_1$ at the weak scale, as in that case the $\\tilde{\\chi}^0_1$ will be the lightest state in a triplet of near mass-degenerate Higgsinos (two neutralinos and one chargino).\\footnote{In the general MSSM it is also possible to have chargino co-annihilation between a pair of wino-dominated $\\tilde{\\chi}^0_1$ and $\\tilde{\\chi}^\\pm_1$, when $|M_2| < |M_1|, |\\mu|$. However, this mechanism is not available in the models discussed here, as the GUT-scale relation $M_1 = M_2 = M_3 \\equiv m_{1\/2}$ leads to $M_2 \\sim 2 M_1$ at the weak scale.}\nIn the \\textsf{CMSSM}\\xspace, the MSSM Higgsino mass parameter $\\mu$ is strongly linked to $m_0$ via the conditions for EWSB; reducing $m_0$ effectively increases $\\mu$. The bino mass parameter $M_1$ is on the other hand controlled by $m_{1\/2}$ via the GUT-scale relation $M_1 = M_2 = M_3 \\equiv m_{1\/2}$. For a fixed value of $m_{1\/2}$, lowering $m_0$ therefore eventually leads to $M_1 \\ll |\\mu|$, resulting in a bino-dominated $\\tilde{\\chi}^0_1$ significantly lower in mass than the Higgsino-dominated neutralinos\/chargino. In the \\textsf{NUHM1}\\xspace, on the other hand, the $\\mu$ parameter is mostly controlled by $m_H$. This allows for $|\\mu| < M_1$, and thus chargino co-annihilation, also in the low-$m_0$ region.\n\n\\begin{figure}[t]\n \\centering\n \\includegraphics[width = 0.49\\textwidth]{figures\/susy\/CMSSM_105_201_like2D_post}\n \\includegraphics[width = 0.49\\textwidth]{figures\/susy\/CMSSM_105_201_ColourMechanism_post}\\\\\n \\includegraphics[height=3.1mm]{figures\/susy\/rdcolours3.pdf}\n \\caption{\n Profile likelihood in the $(m_{\\tilde{\\chi}^0_1},\\Omega_\\chi h^2)$ plane of the \\textsf{CMSSM}\\xspace (\\textit{left}), and the mechanisms that bring the predicted relic density close to or below the measured value (\\textit{right}). The stars show the best-fit points, while the white contours outline the $1\\sigma$ and $2\\sigma$ regions. From~\\cite{CMSSM}.\n }\n \\label{fig:CMSSM_mN1_oh2}\n\\end{figure}\n\nAs mentioned above, in these fits the observed DM relic density is only imposed as an upper bound, to leave open the possibility for non-MSSM contributions in the observed DM density. While this choice broadens the allowed parameter space, it is worth noting that the parameter regions that fully explain the relic density can have equally high likelihoods as those with a lower predicted relic density. This can be seen in the left panel of Fig.~\\ref{fig:CMSSM_mN1_oh2}, which shows the profile likelihood in the \\textsf{CMSSM}\\xspace plane of the neutralino mass $m_{\\tilde{\\chi}^0_1}$ and the predicted relic density $\\Omega_{\\chi} h^2$. For most $m_{\\tilde{\\chi}^0_1}$ values there is little variation in the profile likelihood when moving up to a point where the prediction saturates the observed value (dashed purple line).\n\nThe right-hand panel in Fig.~\\ref{fig:CMSSM_mN1_oh2} shows that, in the \\textsf{CMSSM}\\xspace, the lowest predicted neutralino masses are found within the stop and chargino co-annihilation regions, extending down to $m_{\\tilde{\\chi}^0_1} \\sim 250$\\,GeV. In the \\textsf{NUHM1}\\xspace and \\textsf{NUHM2}\\xspace, the chargino co-annihilation and stau co-annihilation regions extend further down, to $m_{\\tilde{\\chi}^0_1} \\sim 150$\\,GeV. The chargino co-coannihilation region in Fig.~\\ref{fig:CMSSM_mN1_oh2} also illustrates the well-known result that a dominantly Higgsino $\\tilde{\\chi}^0_1$ produces the entire observed relic density when $m_{\\tilde{\\chi}^0_1} \\sim 1$\\,TeV. Moving along the observed relic density towards higher neutralino masses, additional contributions from resonant $A\/H$-funnel annihilations become more and more important.\n\n\\begin{figure}[t]\n \\centering\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/CMSSM_105_217_ColourMechanism_post_wExps.pdf}\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/NUHM1_105_217_ColourMechanism_post_wExps.pdf}\\\\\n \\includegraphics[height=3.1mm]{figures\/susy\/rdcolours4.pdf}\n \\caption{\n The $2\\sigma$ preferred regions in the plane of the spin-independent neutralino-proton cross-section versus the neutralino mass for the \\textsf{CMSSM}\\xspace (\\textit{left}) and the \\textsf{NUHM1}\\xspace (\\textit{right}), coloured according to the mechanism(s) that limit the predicted DM relic density. The pink lines show the observed 90\\% CL exclusion limit from LUX~\\cite{LUXrun2} and projected limits for XENON1T (two tonne-years of exposure), XENONnT\/LZ (20 tonne-years of exposure)~\\cite{XENONnTLZ} and DARWIN (200 tonne-years of exposure)~\\cite{DARWIN}. The $1\\sigma$ and $2\\sigma$ regions are shown as white contours; best-fit points are marked by stars. From~\\cite{CMSSM}.\n }\n \\label{fig:CMSSM_NUHM1_direct_detection}\n\\end{figure}\n\nDirect detection DM searches seem the most promising experimental probe for the SUSY scenarios preferred in these fits. In Fig.\\ \\ref{fig:CMSSM_NUHM1_direct_detection} the preferred \\textsf{CMSSM}\\xspace (left) and \\textsf{NUHM1}\\xspace regions are shown in the plane of the lightest neutralino mass versus the spin-independent neutralino-proton cross-section. The predicted cross-section is scaled by the fraction $f$ of the full DM relic density that the given parameter point attributes to neutralinos. The solid pink line shows the $90\\%$\\,CL exclusion limit from the LUX 2016 result~\\cite{LUXrun2}, which was included as a likelihood component in these fits. The dashed and dotted lines show projected $90\\%$\\,CL limits for the XENON and DARWIN experiments~\\cite{XENONnTLZ, DARWIN}. While the stop co-annihilation regions will largely remain out of reach, as will much of the stau co-annihilation region in the \\textsf{NUHM1}\\xspace, both the chargino co-annihilation and the $A\/H$ funnel regions can be fully probed in future direct detection searches.\n\nFinally, we note that the \\textsf{CMSSM}\\xspace, \\textsf{NUHM1}\\xspace and \\textsf{NUHM2}\\xspace fit results in~\\cite{CMSSM} indicate that these models no longer hold much promise for resolving the observed discrepancy in the muon anomalous magnetic moment. The strong constraints on the low-mass parameter space -- in particular from LHC sparticle searches, DM direct detection and the LHC Higgs measurements -- push the fits towards heavier sfermion and electroweakino spectra, thus diminishing the possible SUSY contribution to the muon $(g-2)$.\n\n\n\n\\subsubsection{Results for the \\textsf{MSSM7}\\xspace}\n\\label{sec:MSSM7}\n\nWe now move on to the weak-scale parameterisations of the MSSM, starting with the \\textsf{GAMBIT}\\xspace analysis of the \\textsf{MSSM7}\\xspace in~\\cite{MSSM}. Here the free parameters are the wino mass parameter, $M_2$; the $(3,3)$ elements of the $\\mathbf{A}_u$ and $\\mathbf{A}_d$ MSSM trilinear coupling matrices, $(\\mathbf{A}_u)_{33} \\equiv A_{u_3}$ and $(\\mathbf{A}_d)_{33} \\equiv A_{d_3}$ (the other trilinear couplings are set to 0); the soft-breaking Higgs mass parameters, $m_{H_u}^2$ and $m_{H_u}^2$; a common parameter $m_{\\tilde{f}}^2$ for the sfermion soft-breaking mass parameters; and the ratio of the Higgs vacuum expectation values, $v_u\/v_d \\equiv \\tan\\beta$. All the parameters are defined at the scale $Q = 1$\\,TeV, except $\\tan\\beta$, which is defined at $Q = m_Z$.\n\nWhile this model is a weak-scale MSSM parameterisation, the GUT-inspired relation\n\\begin{align}\n\\frac{3}{5}\\cos^2\\theta_\\mathrm{W}M_1 = \\sin^2\\theta_\\mathrm{W}M_2 = \\frac{\\alpha}{\\alpha_\\mathrm{s}}M_3,\n\\label{eq:GUT_relation}\n\\end{align}\nis imposed to limit the dimensionality of the parameter space. Equation~\\ref{eq:GUT_relation} represents an expected weak-scale relation between $M_1$, $M_2$ and $M_3$ if they originate from a common GUT-scale parameter, like $m_{1\/2}$ in the \\textsf{CMSSM}\\xspace.\n\nAs in the GUT-scale models, the Higgsino mass parameter $\\mu$ is determined from the input parameters -- most importantly $m_{H_u}^2$ and $m_{H_u}^2$ -- and the requirements for EWSB. Since Eq.\\ \\ref{eq:GUT_relation} implies that $|M_1| < |M_2|$, we again have three \\textit{a priori} possibilities for the composition of the neutralino state: dominantly bino ($|M_1| < |\\mu|$), dominantly Higgsino ($|\\mu| < |M_1|$), or a bino-Higgsino mixture ($|M_1| \\sim |\\mu|$).\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSM7_111_112_like2D.pdf}\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSM7_111_112_ColourMechanism_post.pdf}\\\\\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSM7_105_201_like2D.pdf}\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSM7_105_201_ColourMechanism_post.pdf}\\\\\n \\includegraphics[height=3.1mm]{figures\/susy\/MSSM7_rdcolours5.pdf}\n \\caption{\n Profile likelihoods in the $(\\mu,M_1)$ plane (\\textit{top left}) and the $(m_{\\tilde{\\chi}^0_1},\\Omega_\\chi h^2)$ plane of the \\textsf{MSSM7}\\xspace. The right-hand panels show the $2\\sigma$ preferred parameter regions coloured according to which mechanism(s) contribute to limit the relic density. The stars mark the best-fit points, while the white contours show the $1\\sigma$ and $2\\sigma$ preferred regions. From~\\cite{MSSM}.\n }\n \\label{fig:MSSM7}\n\\end{figure}\n\nThe global fit analysis in~\\cite{MSSM} finds that all these three neutralino scenarios are allowed within the $2\\sigma$ preferred parameter space of the \\textsf{MSSM7}\\xspace. This can be seen in the top panels of Fig.\\ \\ref{fig:MSSM7}, showing the profile likelihood in the $(\\mu, M_1)$ plane (left) and the active mechanisms that bring the relic density close to or below the observed value (right). In the $\\mu < |M_1|$ regions of the plane, corresponding to a mostly Higgsino $\\tilde{\\chi}^0_1$, the chargino co-annihilation and $A\/H$ funnel mechanisms dominate. Moving towards larger $\\mu$ we enter the bino-Higgsino mixture scenario at $\\mu \\sim |M_1|$, before reaching the bino-$\\tilde{\\chi}^0_1$ scenario at $\\mu > |M_1|$. Here the chargino co-annihilation mechanism is no longer relevant, so an acceptable relic density must be achieved either through efficient $A\/H$ funnel annihilations, co-annihilations with the lightest stop or sbottom, or a combination of these mechanisms.\\footnote{The lack of a stau co-annihilation region in the \\textsf{MSSM7}\\xspace is related to the assumption of a common sfermion mass parameter defined at the low scale of $Q=1$\\,TeV. The differences in sfermion masses are then mostly determined by the amount of L\/R mixing in the sfermion mass matrices, rather than RGE running of mass parameters. Since the L\/R mixing terms for both up-type and down-type sfermions are proportional to the corresponding Yukawa couplings, the light stop ends up being the lightest sfermion across much of parameter space, and the light sbottom is always lighter than then light stau.}\n\n\nThe overall best-fit point in the \\textsf{MSSM7}\\xspace is found in the chargino co-annihilation region, with $m_{\\tilde{\\chi}^0_2} \\approx m_{\\tilde{\\chi}^\\pm_1} \\approx m_{\\tilde{\\chi}^0_1} \\approx 260$\\,GeV. As can be seen in the lower panels of Fig.\\ \\ref{fig:MSSM7}, the predicted neutralino relic density at this point can only explain around 10\\% of the observed DM relic density. However, with only slightly heavier neutralino masses there are \\textsf{MSSM7}\\xspace scenarios that achieve close to the same likelihood values -- well within the $1\\sigma$ region -- and account for the full relic density. These are scenarios with a mostly bino $\\tilde{\\chi}^0_1$ and efficient $\\tilde{\\chi}^0_1$--$\\tilde{\\chi}^0_1$ annihilations through the $A\/H$ funnel.\n\nThe cutoff of this $A\/H$ funnel region at $m_{\\tilde{\\chi}^0_1} \\sim 250$\\,GeV, corresponding to $m_{A\/H} \\sim 500$\\,GeV, is due to several independent likelihood contributions that penalize the lower-mass scenarios. In particular, the constraint on BSM contributions to $BR(B \\rightarrow X_s \\gamma)$ plays an important role here, as the $A^0$ mass is closely related to the $H^\\pm$ mass, and a light charged Higgs will induce sizable SUSY contributions to this decay. Further important constraints on this region come from the LHC Higgs measurements, and also from LHC gluino searches, as the gluino mass parameter $M_3$ is connected to $M_1$ via Eq.\\ \\ref{eq:GUT_relation}, giving $M_3 \\sim 5 M_1$.\n\n\\begin{figure}[t]\n \\centering\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSM7_105_217_like2D.pdf}%\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSM7_105_217_ColourMechanism_post.pdf}\\\\\n \\includegraphics[height=3.1mm]{figures\/susy\/MSSM7_rdcolours5.pdf}\n \\caption{\n\tProfile likelihood in the plane of the neutralino mass versus the spin-independent neutralino-proton cross-section in the \\textsf{MSSM7}\\xspace (left), and the relic density mechanisms that are active in different parts of the $2\\sigma$ region (right). The predicted neutralino-proton cross-section is rescaled at each point by the fraction $f$ of the observed DM relic density that the neutralino relic prediction accounts for. 90\\% CL exclusion limits are shown for the full LUX exposure~\\cite{LUXrun2} and the projected reach for for XENON1T (two tonne-years of exposure), XENONnT\/LZ (20 tonne-years of exposure)~\\cite{XENONnTLZ} and DARWIN (200 tonne-years of exposure)~\\cite{DARWIN}. The $1\\sigma$ and $2\\sigma$ regions are outlined by white contours. The stars mark the best-fit points. From~\\cite{MSSM}.\n }\n \\label{fig:MSSM7_direct_detection}\n\\end{figure}\n\nWe note that even the $h\/Z$ funnel mechanisms are present within the $2\\sigma$ parameter regions, for $m_{\\tilde{\\chi}^0_1} \\approx 45$\\,GeV and $m_{\\tilde{\\chi}^0_1} \\approx 62$\\,GeV. However, the allowed scenarios in this low-$m_{\\tilde{\\chi}^0_1}$ region have an almost pure Higgsino $\\tilde{\\chi}^0_1$ anyway, so this alone ensures a predicted relic density far below the observed value, also explaining how the otherwise strong constraints from DM direct detection are avoided.\n\nAs for the GUT-scale models discussed in the previous section, direct DM searches seem the most promising probe of the \\textsf{MSSM7}\\xspace scenarios preferred by this fit. Figure \\ref{fig:MSSM7_direct_detection} shows the profile likelihood (left) and the active DM mechanisms (right) across the plane of the neutralino mass and the spin-independent neutralino-proton cross-section. We see that future direct detection experiments will explore not only the full chargino co-annihilation region, but almost the entire $1\\sigma$ region preferred in the \\textsf{GAMBIT}\\xspace fit.\n\nConcerning the muon $(g-2)$ discrepancy, the fit in~\\cite{MSSM} shows that there is little hope that the \\textsf{MSSM7}\\xspace can provide an explanation. This is not particularly surprising: because the model dimensionality is kept low, relating all sfermion mass parameters to the common $m_{\\tilde{f}}^2$ parameter at the weak scale, it is impossible to get sufficiently light smuons and muon sneutrinos without simultaneously causing significant tension with other observables such as LHC squark searches.\n\n\n\n\\subsubsection{Results for the \\textsf{EWMSSM}\\xspace}\n\\label{sec:EWMSSM}\n\nCurrent SUSY searches by the ATLAS and CMS experiments at the LHC are usually optimised and interpreted assuming a simplified model. These models typically include only two or three different sparticles and assume 100\\% of decays occur to the signal processes. Such theory simplifications are a necessary compromise given the level of detail and complexity in experimental searches. Nevertheless it leaves open an important question: what impact do the results from ATLAS and CMS SUSY searches have on the parameter space of more realistic models like the MSSM?\n\nThe \\textsf{GAMBIT}\\xspace analysis in~\\cite{EWMSSM} takes on this question in the context of LHC searches for neutralinos and charginos. The canonical simplified model for these searches is one that assumes production of a purely wino $\\tilde{\\chi}^0_2 \\tilde{\\chi}^\\pm_1$ pair, with subsequent decays to a purely bino $\\tilde{\\chi}^0_1$ via $\\tilde{\\chi}^0_2 \\rightarrow Z \\tilde{\\chi}^0_1$ and $\\tilde{\\chi}^\\pm_1 \\rightarrow W^\\pm \\tilde{\\chi}^0_1$. This gives motivation for a search for events with leptons, jets and missing energy (see e.g.\\ \\cite{Aaboud:2018jiw, Sirunyan:2017lae}). The \\textsf{GAMBIT}\\xspace study assumes a phenomenologically far richer model, referred to as the \\textsf{EWMSSM}\\xspace. This is the effective theory obtained when assuming that all sparticles except the MSSM electroweakinos are too heavy to affect current collider searches. The \\textsf{EWMSSM}\\xspace is thus a model with six sparticles -- four neutralinos and two charginos -- controlled by only four free MSSM parameters: $M_1$, $M_2$, $\\mu$ and $\\tan\\beta$. Loosely speaking, the bino soft-mass $M_1$ controls the mass of one neutralino, the wino soft-mass $M_2$ controls the masses of one neutralino and one chargino, and the Higgsino mass parameter $\\mu$ sets the masses of two neutralinos and one chargino.\n\nIn contrast to the global fits discussed in the previous two sections, the fit in \\cite{EWMSSM} focuses exclusively on collider constraints. This choice allows the fit to explore the full range of possible collider scenarios in the \\textsf{EWMSSM}\\xspace without further enlarging the model parameter space. Keeping the dimensionality of the parameter space fairly low is of critical importance, due the large computational expense of this fit: for each sampled \\textsf{EWMSSM}\\xspace parameter point, \\textsf{ColliderBit}\\xspace is used to run full Monte Carlo simulations of the relevant ATLAS and CMS searches. While running full simulations at each point in a global fit is always computationally challenging, it is particularly so when simulating electroweakino searches due to the low signal acceptance rates in these searches.\\footnote{For most of the included LHC searches there is no public information on how background estimates are correlated across signal regions. In these cases the single signal region with the best expected sensitivity must be identified at each \\textsf{EWMSSM}\\xspace parameter point in the fit. This adds to the already substantial computational cost, as distinguishing between ``competing'' signal regions often requires higher Monte Carlo statistics than what is needed to get reasonable signal estimates for each individual signal region alone.}\n\nThe analysis in \\cite{EWMSSM} includes \\textsf{ColliderBit}\\xspace simulations of most of the $13$\\,TeV electroweakino searches that were available at the time of the study~\\cite{Aaboud:2018jiw,Aaboud:2018sua,Aaboud:2018htj,Aaboud:2018zeb,CMS:2017fth,Sirunyan:2018iwl,Sirunyan:2017qaj,CMS-PAS-SUS-16-039}. The combined likelihood obtained from these simulations is the main component in the fit likelihood function. The other collider observables going into the total likelihood are a collection of SUSY cross-section limits from LEP and the invisible decay widths of the $Z$ and the $125$\\,GeV Higgs.\n\n\\begin{figure}[t]\n \\centering\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSMEW_155_151_like2D.pdf}\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSMEW_152_155_like2D.pdf}\\\\\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSMEW_152_153_like2D.pdf}\n \\includegraphics[width=0.49\\textwidth]{figures\/susy\/MSSMEW_154_153_like2D.pdf}\n \\caption{\n Profile likelihood in four different \\textsf{EWMSSM}\\xspace mass planes: the $(m_{\\tilde{\\chi}^\\pm_1},m_{\\tilde{\\chi}^0_1})$ plane (top left), the $(m_{\\tilde{\\chi}^0_2},m_{\\tilde{\\chi}^\\pm_1})$ plane (top right), the $(m_{\\tilde{\\chi}^0_2},m_{\\tilde{\\chi}^0_3})$ plane (bottom left), and the $(m_{\\tilde{\\chi}^0_4},m_{\\tilde{\\chi}^0_3})$ plane (bottom right). The white contours show the $1\\sigma$ and $2\\sigma$ preferred regions. The star marks the best-fit point. From~\\cite{EWMSSM}.\n }\n \\label{fig:EWMSSM_mass_planes}\n\\end{figure}\n\nThe main result from \\cite{EWMSSM} is that, when combined, the ATLAS and CMS electroweakino results prefer \\textsf{EWMSSM} scenarios with a distinct pattern of relatively light neutralino and chargino masses (Fig.\\ \\ref{fig:EWMSSM_mass_planes}). The preferred $2\\sigma$ parameter region has all six neutralinos and charginos below $\\sim$$700$\\,GeV, with the lightest neutralino below $\\sim$$200$\\,GeV. The lightest neutralino is always dominantly bino, but it also has a non-negligible wino or Higgsino component. Further, the best-fit parameter region predicts two characteristic $\\gtrsim m_Z$ gaps in the mass spectrum: the first between the mostly bino $\\tilde{\\chi}^0_1$ and the mostly wino (Higgsino) $\\tilde{\\chi}^0_2$\/$\\tilde{\\chi}^\\pm_1$, and the second between $\\tilde{\\chi}^0_2$\/$\\tilde{\\chi}^\\pm_1$ and the mostly Higgsino (wino) $\\tilde{\\chi}^0_4$\/$\\tilde{\\chi}^\\pm_2$.\\footnote{In the preferred scenarios, $\\tilde{\\chi}^0_3$ is always mostly Higgsino and thus fairly close in mass to the other Higgsino-dominated states, i.e.\\ either $\\tilde{\\chi}^0_2$\/$\\tilde{\\chi}^\\pm_1$ or $\\tilde{\\chi}^0_4$\/$\\tilde{\\chi}^\\pm_2$.}\n\nAt first sight this result may seem surprising. None of the included ATLAS and CMS searches have seen a convincing SUSY signal, yet when combined they prefer the low-mass region over the decoupling region, where all \\textsf{EWMSSM}\\xspace collider predictions would align with SM expectations. The reason is that the \\textsf{EWMSSM}\\xspace is able to simultaneously fit a pattern of small excesses across several of the simulated LHC searches, while at the same time avoiding generating too much tension with the other searches. The excesses that mostly drive this result come from searches for 2-, 3-, and 4-lepton final states in ATLAS \\cite{Aaboud:2018jiw,Aaboud:2018sua,Aaboud:2018zeb}, specifically in signal regions that target leptons from on-shell $Z$ and $W$ decays. This explains the preference in the fit for electroweakino mass spectra with two $\\gtrsim m_Z$ mass gaps.\n\n\\begin{figure}\n \\centering\n %\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_155_151_obs2D_605.pdf}\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_152_153_obs2D_605.pdf}\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_154_153_obs2D_605.pdf}\n %\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_155_151_obs2D_607.pdf}\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_152_153_obs2D_607.pdf}\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_154_153_obs2D_607.pdf}\n %\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_155_151_obs2D_608.pdf}\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_152_153_obs2D_608.pdf}\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_154_153_obs2D_608.pdf}\n %\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_155_151_obs2D_610.pdf}\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_152_153_obs2D_610.pdf}\n \\includegraphics[width=0.32\\textwidth]{figures\/susy\/MSSMEW_154_153_obs2D_610.pdf}\n %\n \\caption{\n Contributions to the total fit likelihood from the ATLAS searches in Ref.\\ \\cite{Aaboud:2018zeb} (top), Ref.\\ \\cite{Aaboud:2018jiw} (second and third rows), and Ref.\\ \\cite{Aaboud:2018sua} (bottom), shown across the full $3\\sigma$ regions in the $(m_{\\tilde{\\chi}^\\pm_1}, m_{\\tilde{\\chi}^0_1})$ plane (left), the $(m_{\\tilde{\\chi}^0_2}, m_{\\tilde{\\chi}^0_3})$ plane (middle), and the $(m_{\\tilde{\\chi}^0_4}, m_{\\tilde{\\chi}^0_3})$ plane (right). In the blue regions a non-zero signal prediction in the given search improves the overall fit, while in red regions the signal prediction worsens the fit. In the white regions the given search is not sensitive. The orange contours outline the $1\\sigma$, $2\\sigma$ and $3\\sigma$ regions preferred in the fit. The white star marks the best-fit point. From~\\cite{EWMSSM}.\n }\n \\label{fig:EWMSSM_mass_planes_per_analysis}\n\\end{figure}\n\nTo understand the interplay between the analyses contributing to the excess, we can look at their individual likelihood contributions across the combined best-fit surface. This is done in Fig.\\ \\ref{fig:EWMSSM_mass_planes_per_analysis}, where the contributions from four ATLAS results are displayed across the preferred $3\\sigma$ regions in three different mass planes. When reading these plots it is important to keep in mind that the plotted points are those parameter samples picked out by profiling the \\textit{total} likelihood. The sharp changes in analysis likelihood seen in some plots are due to abrupt changes in what scenarios are picked out by this profiling, which again changes which signal region is selected to set the analysis likelihood value.\n\nOne example of the interplay between analyses is seen by comparing the middle panels on the first and third rows. The first of these show the likelihood contribution from an ATLAS search for 4-lepton final states, with the leptons coming from two $Z$ bosons~\\cite{Aaboud:2018zeb}. We see that fitting a 4-lepton excess in the \\textsf{EWMSSM}\\xspace relies on having non-negligible production of $\\tilde{\\chi}^0_3$, as this allows for signal leptons from the decays $\\tilde{\\chi}^0_3 \\rightarrow Z \\tilde{\\chi}^0_{1,2}$. The second of these panels is for an ATLAS search for 3-lepton final states~\\cite{Aaboud:2018jiw}, designed to target $\\tilde{\\chi}^0_2 \\tilde{\\chi}^\\pm_1$ production. For a given $m_{\\tilde{\\chi}^0_2}$, reducing $m_{\\tilde{\\chi}^0_3}$ to $\\lesssim 600$\\,GeV (as preferred by the 4-lepton search) also improves the fit to this 3-lepton search, which for high $m_{\\tilde{\\chi}^0_3}$ sees some tension with the data. At lower $m_{\\tilde{\\chi}^0_3}$ production processes with $\\tilde{\\chi}^0_3$ come into play, involving more complicated event topologies. At the same time the production cross-section for the $\\tilde{\\chi}^0_2 \\tilde{\\chi}^\\pm_1$ pair is reduced somewhat, due to a higher Higgsino component. The combined effect is a change in which 3-lepton signal region is identified as having the best expected sensitivity.\n\nThe combined excess in the $13$\\,TeV searches is estimated in~\\cite{EWMSSM} to have a local significance of $3.3\\sigma$. The impact of $8$\\,TeV LHC results on the preferred low-mass scenarios is investigated by post-processing all parameter samples in the $1\\sigma$ region with simulations of relevant ATLAS and CMS electroweakino searches at $8$\\,TeV \\cite{Aad:2015jqa,ATLAS:2LEPEW_20invfb,ATLAS:3LEPEW_20invfb,CMS:3LEPEW_20invfb}. The result is an upwards shift in the best-fit mass spectrum, by $\\sim20$\\,GeV in all masses, and a small reduction of the estimated significance of the excess, to $2.9\\sigma$.\n\nWe also note that even though the \\textsf{EWMSSM}\\xspace fit did not include DM constraints, parts of the preferred parameter space do give acceptable relic density predictions while avoiding exclusion from current direct and indirect DM searches. This is possible for scenarios with $m_{\\tilde{\\chi}^0_1}$ close to $m_Z\/2$ or $m_h\/2$, where resonant annihilations via the $Z\/h$ funnel can bring the predicted relic density close to or below the observed value.\n\nWhile the small excess seen in the \\textsf{EWMSSM}\\xspace fit is quite possibly due to background fluctuations, the fit demonstrates two important points. First, that LHC constraints on light SUSY can be significantly weaker in realistic SUSY such as the MSSM than in simplified models.\\footnote{While not discussed here, the analysis in~\\cite{EWMSSM} shows that for every mass hypothesis in the $(m_{\\tilde{\\chi}^0_2}, m_{\\tilde{\\chi}^0_1})$ plane -- not just for points in the best-fit region -- there is a point in the \\textsf{EWMSSM}\\xspace parameter space that fits the \\textit{combined} collider results at least as well as the SM expectation.} Second, that proper statistical combinations of collider searches can be a powerful tool to uncover suggestive patterns in BSM parameter spaces.\n\n\n\n\\subsection{Higgs Portal models for dark matter}\n\nNo definitive evidence has yet been uncovered for non-gravitational interactions of DM with the SM. At some level however, such interactions must be inevitably generated by effective operators connecting Lorentz-invariant, gauge singlet combinations of SM particles to equivalently symmetric combinations of DM fields. The lowest-dimension such operator in the SM is the Higgs bilinear $H^\\dagger H$. Depending on the spin and gauge representation of a DM candidate $X$, the lowest-order Lorentz- and gauge-invariant DM operator may be either the bilinear $X^\\dagger X$, or a lone DM field. Operators linear in $X$ are only consistent if $X$ is itself a Lorentz invariant (i.e.\\ a scalar), and a gauge singlet. If it is to be a viable DM candidate however, $X$ must be stable on cosmological timescales. The most straightforward way to achieve this is for $X$ to hold a different charge to SM particles under some new unbroken (typically discrete) symmetry. This has the effect of forbidding terms linear in $X$, preventing the field from decaying.\n\nThe lowest-order operator connecting $X$ to the SM guaranteed to exist at some level is therefore the so-called `Higgs portal' operator $X^\\dagger X H^\\dagger H$. Following electroweak symmetry breaking, this operator gives rise to a mass term for $X$ proportional to $v_0^2$ (with $v_0$ the vacuum expectation value of the Higgs field), a Higgs-DM-DM vertex proportional to $v_0$, and a direct four-particle vertex between two Higgses and two DM particles. The new 3-particle and 4-particle interactions of $X$ with the Higgs boson lead to DM annihilation (enabling thermal production and possible indirect detection), spin-independent DM-nucleon scattering (leading to possible direct detection), DM production at colliders (with the possibility for signals in e.g.\\ monojet searches), and invisible decays of the Higgs to two DM particles when $m_X < m_h\/2$.\n\nDepending on the Lorentz representation of DM, $X^2 H^2$ may be a fully renormalisable dimension 4 operator (if $X$ is a scalar), an effective dimension 4 operator (if $X$ is a vector), or an effective dimension 5 operator (if $X$ is a fermion). All three of these cases have been considered in detail in the literature, with a particular focus on models where $X$ is itself a gauge singlet and the $X^2H^2$ term is therefore the sole link between DM and the SM. The most commonly studied cases have been the $\\mathbb{Z}_2$-symmetric scalar [\\citenum{SilveiraZee,McDonald94,Burgess01,Davoudiasl:2004be,Goudelis09,Yaguna09,Profumo2010a,Andreas:2010dz,Arina11,Mambrini11, Raidal:2011xk,Mambrini:2011ik,He:2011de,Drozd:2011aa,Okada:2012cc,Cheung:2012xb,Okada:2013bna,Cline:2013gha,Chacko:2013lna, Endo:2014cca,Craig:2014lda, Feng15,Duerr15,arXiv:1510.06165,Duerr16,He:2016mls,Han:2016gyy,Dupuis:2016fda,Cuoco:2016jqt,Binder:2017rgn,Ghorbani:2018yfr,Chiang:2018gsn,Stocker:2018avm,Hardy:2018bph,Bernal:2018kcw,Glioti:2018roy,Urbano:2014hda,Escudero:2016gzx,Kanemura:2011nm,Djouadi:2011aa,Djouadi:2012zc,Bishara:2015cha,Ko:2016xwd,Beniwal:2015sdl,Kamon:2017yfx,Dutta:2017sod,Dick:2018lqx}; \\textsf{GAMBIT}\\xspace analyses \\citenum{SSDM,SSDM2}], vector [\\citenum{Djouadi:2011aa,Kanemura:2011nm,Djouadi:2012zc,Bishara:2015cha,Chen:2015dea,DiFranzo:2015nli, Beniwal:2015sdl,Ko:2016xwd,Kamon:2017yfx,Dutta:2017sod,arXiv:1704.05359,Dick:2018lqx,Baek:2014jga}; \\textsf{GAMBIT}\\xspace analysis \\citenum{HP}] and fermionic [\\citenum{Djouadi:2011aa,Kanemura:2011nm,LopezHonorez:2012kv,Djouadi:2012zc,Urbano:2014hda, Baek:2014jga,Bishara:2015cha,Beniwal:2015sdl,Ko:2016xwd,Fedderke:2014wda,Matsumoto:2014rxa,arXiv:1506.04149, arXiv:1506.08805,arXiv:1506.06556, Escudero:2016gzx,Kamon:2017yfx,Dutta:2017sod,Dick:2018lqx,Matsumoto:2018acr}; \\textsf{GAMBIT}\\xspace analysis \\citenum{HP}] variants, along with the $\\mathbb{Z}_3$-symmetric scalar [\\citenum{Belanger2013a,Kang:2017mkl,2017JHEP...10..088B,Hektor:2019ote,Kannike:2019mzk}; \\textsf{GAMBIT}\\xspace analysis \\citenum{SSDM2}].\n\n\n\\subsubsection{$\\mathbb{Z}_2$-symmetric scalar singlet}\n\nThe simplest Higgs portal model for DM, and indeed probably the most minimal of all models for particle DM, is a single, real, gauge-singlet scalar field $S$, protected from decay by a $\\mathbb{Z}_2$ symmetry. The only new renormalisable Lagrangian terms allowed by gauge, Lozentz and $\\mathbb{Z}_2$ symmetry are\n\\begin{equation}\n\\mathcal{L}_{\\mathbb{Z}_2} = \\frac12 \\mu_{\\scriptscriptstyle S}^2 S^2 + \\frac14\\lambda_{\\sss S} S^4 + \\frac12\\lambda_{h\\sss S} S^2|H|^2.\n\\label{L_S}\n\\end{equation}\nThe model is fully specified by the $S$ bare mass $\\mu_{\\scriptscriptstyle S}$, the dimensionless $S$ quartic self-coupling $\\lambda_{\\sss S}$, and the dimensionless Higgs portal coupling $\\lambda_{h\\sss S}$. For the most part, the $S$ quartic coupling has little impact on the phenomenology of the model, as it leads only to DM self-interactions, which are not sufficiently constrained by existing data to place strong limits on $\\lambda_{\\sss S}$. A key exception, however, is the impact of $\\lambda_{\\sss S}$ on the running of gauge couplings under renormalisation group flow, which can have important implications for stability of the electroweak vacuum.\n\nDenoting the physical SM Higgs field by $h$, following electroweak symmetry breaking $H \\rightarrow \\left[0, (v_0+h)\/\\sqrt{2}\\right]^\\text{T}$. This generates new vertices of the form $v_0hS^2$ and $h^2S^2$, and induces a shift to the $S$ bare mass, such that at tree level\n\\begin{equation}\nm_{\\sss S} = \\sqrt{\\mu_{\\scriptscriptstyle S}^2 + \\frac12{\\lambda_{h\\sss S} v_0^2}}.\n\\label{ms}\n\\end{equation}\n\n\\begin{figure}[t]\n\\centering\n\\includegraphics[width=0.6\\columnwidth]{figures\/figure_1.pdf}\n\\caption{Feynman diagrams for annihilation, semi-annihilation, nuclear scattering and Higgs decays in scalar singlet Higgs portal models. $N, f$ and $V$ refer to nucleons, fermions and SM electroweak vector bosons ($Z$ and $W$), respectively. Diagrams are shown for the $\\mathbb{Z}_3$-symmetric case, where DM exists in $S$ and anti-$S$ (i.e. $S^*$) states, but the same diagrams apply in the $\\mathbb{Z}_2$-symmetric case with $S=S^*$, except for semi-annihilation (which is absent in the $\\mathbb{Z}_2$ model). The same diagrams also apply to $\\mathbb{Z}_2$-symmetric vector and fermionic Higgs portal models (with $S$ replaced by the relevant DM particle and semi-annihilation also forbidden by the $\\mathbb{Z}_2$ symmetry). From \\cite{SSDM2}.}\n\\label{fig:diagrams}\n\\end{figure}\n\n\\begin{figure}[t]\n\t\\centering\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_singlet_19_17_like2D_SingletDM_Z2_low_X}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_singlet_19_17_like2D_SingletDM_Z2_full_X}\n \\caption{Profile likelihoods of parameters in the $\\mathbb{Z}_2$-symmetric scalar singlet Higgs portal dark matter model, including constraints from direct and indirect detection, the relic density of dark matter and LHC searches for invisible decays of the Higgs boson, along with various Standard Model, dark matter halo and nuclear uncertainties. \\textit{Left}: the low-mass resonance region. \\textit{Right}: the full mass range. Contours show 1 and 2$\\sigma$ confidence regions, with white corresponding to the main scan (including the 2018 XENON1T direct search \\cite{Aprile:2018dbl}) and grey to a secondary scan using the 2017 XENON1T result \\cite{Aprile:2017iyp}. White stars indicate the location of the best-fit point. From \\protect\\cite{SSDM2}.}\n\t\\label{fig:z2scalar1}\n\\end{figure}\n\n\\begin{figure}[t]\n\t\\centering\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_singlet_19_200_like2D_SingletDM_Z2_full_X}\n \\caption{Results from the same analysis of the $\\mathbb{Z}_2$-symmetric scalar singlet Higgs portal dark matter model as shown in Fig.\\ \\ref{fig:z2scalar1}, but plotted in the plane of the effective spin-independent nuclear scattering cross-section and the scalar mass, in order to compare directly to the sensitivity of direct detection experiments. All models have their effective cross-section defined as $f\\sigma_\\mathrm{SI}$, where $f\\equiv \\Omega_S \/ \\Omega_\\mathrm{DM}$ is the fraction of the relic density constituted by the scalar singlet. Experiments assume $f=1$ when publishing their results. Contours show 1 and 2$\\sigma$ confidence regions, and stars best fits. From \\protect\\cite{SSDM2}.}\n\t\\label{fig:z2scalar2}\n\\end{figure}\n\n\\begin{figure}[t]\n\t\\centering\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_singlet_19_17_like2D_SingletDM_Z2_full_vs_X_cut}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_singlet_19_200_like2D_SingletDM_Z2_full_vs_X_cut}\n \\caption{Regions in the $\\mathbb{Z}_2$-symmetric scalar singlet model that satisfy all experimental constraints, stabilise the electroweak vacuum and remain perturbative up to scales of $10^{15}$\\,GeV. Contours show 1 and 2$\\sigma$ confidence regions, and stars best fits. Grey contours show the allowed regions without the requirements of vacuum stability and perturbativity. From \\protect\\cite{SSDM2}.}\n\t\\label{fig:z2scalar3}\n\\end{figure}\n\nThe interaction with the physical Higgs endows $S$ with essentially all of the classic phenomenology of WIMP DM, via the diagrams shown in Fig.\\ \\ref{fig:diagrams} -- along with the added possibility of Higgs decays $h\\to SS$ where $m_{\\sss S} \\le m_h\/2$. The leading constraints on the model come from searches for gamma rays from dark matter annihilation in dwarf spheroidal galaxies \\cite{LATdwarfP8}, the observed relic density of dark matter \\cite{Planck18cosmo}, direct searches performed by the XENON1T \\cite{Aprile:2018dbl} and PandaX \\cite{Cui:2017nnn} experiments, and searches for invisible Higgs decays at the LHC \\cite{Belanger:2013xza,CMS-PAS-HIG-17-023}.\n\nThe resulting preferred regions of parameter space are shown in Fig.\\ \\ref{fig:z2scalar1}. These results explicitly allow models where $S$ is only a fraction of the observed DM, and include a fully self-consistent rescaling of the predicted signals at direct and indirect searches according to the fraction $f \\le 1$ of DM constituted by $S$ at each point in the parameter space. The allowed parameter space splits into three regions: one at high masses where direct detection loses sensitivity, a second at intermediate mass where the 4-boson vertex boosts the annihilation cross-section and depletes the relic density, and another at and immediately below $m_{\\sss S} = m_h\/2$, where $S$ annihilates highly efficiently via an $s$-channel resonance mediated by the Higgs, depleting the relic density to below the observed value even for very small values of $\\lambda_{h\\sss S}$.\n\nThe Higgs invisible width constraint rules out large couplings $\\lambda_{h\\sss S}$ at singlet masses below the resonance. The thermal relic density of $S$ provides the lower limit of the low-mass and high-mass allowed regions. Indirect detection plays the leading role only on the high-mass edge of the resonance, where thermal effects in the early Universe push annihilation slightly off resonance but late-time annihilation remains strongly boosted. Direct detection plays a significant role throughout the parameter space, as can be seen in Fig.\\ \\ref{fig:z2scalar2}. Except for the very bottom of the resonance region, the entirety of the model will soon be probed by direct detection.\n\nGamma-ray lines do not provide any meaningful constraint, as the partial annihilation cross-section for $SS\\to\\gamma\\gamma$ is only appreciable in parts of the parameter space where the relic density is significantly suppressed. Likewise, monojet searches only constrain very large values of $\\lambda_{h\\sss S}$ already excluded by other constraints or expected to lead to new strong dynamics. Indeed, both these points also apply to all other Higgs portal models that we discuss in this review.\n\nGiven that the Higgs portal operator is not just an effective interaction, but a fully renormalisable operator in this model, it is also important to consider the UV behaviour of the theory. Due to the observed values of the top and Higgs masses, the SM posesses a second minimum in its scalar potential at $\\gtrsim \\mathcal{O}(10^{15})$\\,GeV, causing the low-scale vacuum in which we reside to be metastable. Adding an additional scalar to the SM impacts the running of the Higgs quartic coupling, raising its value at high scales. This can prevent the quartic coupling from running negative, and make the low-scale minimum a global rather than a local one. The catch is that $\\lambda_{\\sss S}$ must be relatively large in order to achieve this effect. Fig.\\ \\ref{fig:z2scalar3} shows the parts of the parameter space, consistent with all experimenal constraints, where $\\lambda_{\\sss S}$ can be pushed high enough to stabilise the SM vacuum, but without pushing any of the couplings non-perturbative below a scale of $10^{15}$\\,GeV. Clearly, the $\\mathbb{Z}_2$-symmetric scalar singlet can solve the vacuum stability problem without introducing new strong dynamics, and satistfy all experimental constraints, but only in a region around \\mbox{$m_{\\sss S} = 1$--2\\,TeV} and $\\sigma_\\mathrm{SI} \\sim 10^{-45}$\\,cm$^2$. Curiously, this is also in the region consistent with the (admittedly very small) excess seen in the most recent XENON1T results \\cite{Aprile:2018dbl}. In any case, this hypothessis will clearly be tested very quickly in the upcoming runs of the LZ and XENONnT \\cite{Akerib:2015cja,XENONnTLZ} experiments.\n\nThe results in Figs.\\ \\ref{fig:z2scalar1}--\\ref{fig:z2scalar3} are based on profile likleihood analyses, and illustrate what is possible in each parameter plane, were one able to freely vary the other parameters of the theory (including nuisance parameters) in order to achieve the best possible fit to all available data. If one instead carries out a Bayesian analysis, looking instead at the posterior probability density for these parameters, a different picture emerges. In this case, parameter combinations become more likely if they can provide a good fit for a broader range of values of the other parameters of the theory, i.e.\\ if they can fit the data with less fine tuning. In this case, the low-mass resonance region is strongly disfavoured, as `hitting' the resonance and avoiding the relic density constraint for a given value of $m_{\\sss S}$ requires some fine-tuning of various SM nuisance parameters such as $m_h$; the same is true to a lesser extent for the intermediate-mass region as well. We therefore see that from a Bayesian perspective, the region where the singlet model stabilises the SM vacuum is in fact favoured over the other regions of the theory, even before considering the implications for vacuum stability.\n\n\\begin{figure}[t]\n\t\\centering\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_singlet_20_17_like2D_SingletDM_Z3_full_X}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_singlet_20_19_like2D_SingletDM_Z3_full_X}\n \\caption{Profile likelihoods of parameters in the $\\mathbb{Z}_3$-symmetric scalar singlet Higgs portal dark matter model, including constraints from direct and indirect detection, the relic density of dark matter and LHC searches for invisible decays of the Higgs boson, along with various Standard Model, dark matter halo and nuclear uncertainties. Contours show 1 and 2$\\sigma$ confidence regions, with white corresponding to the main scan (including the 2018 XENON1T direct search \\cite{Aprile:2018dbl}) and grey to a secondary scan using the 2017 XENON1T result \\cite{Aprile:2017iyp}. White stars indicate the location of the best-fit point. From \\protect\\cite{SSDM2}.}\n\t\\label{fig:z3scalar1}\n\\end{figure}\n\n\\begin{figure}[t]\n\t\\centering\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_singlet_20_17_obs2D_60_SingletDM_Z3_full_X}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_singlet_20_19_obs2D_60_SingletDM_Z3_full_X}\n \\caption{Results from the same analysis of the $\\mathbb{Z}_3$-symmetric scalar singlet Higgs portal dark matter model as shown in Fig.\\ \\ref{fig:z3scalar1}, but shaded according to the semi-annihilation fraction $\\alpha$ (Eq.\\ \\protect\\ref{eqn:sa_fraction}). From \\protect\\cite{SSDM2}.}\n\t\\label{fig:z3scalar2}\n\\end{figure}\n\n\\subsubsection{$\\mathbb{Z}_3$-symmetric scalar singlet}\n\nIn contrast to the self-adjoint $\\mathbb{Z}_2$-symmetric scalar singlet, a $\\mathbb{Z}_3$ symmetry leads to a complex scaler DM candidate, with both DM ($S$) and anti-DM ($S^*$) states contributing to the relic density. This symmetry also allows an additional cubic term in the Lagrangian,\n\\begin{equation}\n\\mathcal{L}_{\\mathbb{Z}_3} = \\mu_{\\scriptscriptstyle S}^2 S^\\dagger S + \\lambda_{\\sss S} (S^\\dagger S)^2 + \\frac{\\mu_3}{2}(S^{\\dagger 3}+S^3) + \\lambda_{h\\sss S} S^\\dagger S|H|^2,\n\\end{equation}\nwhere we have introduced the new dimension-1 $S$ cubic coupling $\\mu_3$. This new coupling allows for so-called semi-annihilation processes $SS\\to S^*h$ and $S^*S^*\\to Sh$, shown in Fig.\\ \\ref{fig:diagrams}.\n\nCompared to the $\\mathbb{Z}_2$-symmetric model, semi-annihilation is able to deplete the relic density of DM at intermediate masses and open up an entirely new region of viable parameter space. This is shown in terms of the profile likelihood in Fig.\\ \\ref{fig:z3scalar1}, and highlighted in terms of the semi-annihilation fraction\n\\begin{equation}\n\\alpha=\\frac{1}{2}\\frac{\\langle\\sigma v_\\mathrm{rel}\\rangle_{SS\\rightarrow hS}}{\\langle\\sigma v_\\mathrm{rel}\\rangle+\\frac{1}{2}\\langle\\sigma v_\\mathrm{rel}\\rangle_{SS\\rightarrow hS}},\\label{eqn:sa_fraction}\n\\end{equation}\nin Fig.\\ \\ref{fig:z3scalar2}. Here $\\langle\\sigma v_\\mathrm{rel}\\rangle$ is the thermally averaged (semi-)annihilation cross-section weighted by the relative velocity between annihilating particles.\n\nThe vacuum structure of the theory is also more complicated than that of the $\\mathbb{Z}_2$-symmetric model, as regions where $\\mu_3 \\geq 2\\sqrt{\\lambda_{\\sss S}}m_{\\sss S}$ or $\\mu_{\\scriptscriptstyle S}^2<0$ and $\\lambda_{h\\sss S}$ is large can possess a second, $\\mathbb{Z}_3$-breaking minimum. The results shown in Figs.\\ \\ref{fig:z3scalar1} and \\ref{fig:z3scalar2} avoid these regions, demanding that $S$ does not itself obtain a VEV, and that the potential remains bounded from below.\n\nLike the $\\mathbb{Z}_2$-symmetric variant, the $\\mathbb{Z}_3$-symmetric model can in principle completely stabilise the SM vacuum. However, because of the various factors of 2 introduced relative to the $\\mathbb{Z}_2$ case, by virtue of DM not being self-adjoint, the region where this is possible is in fact in strong tension with the results from both XENON1T \\cite{Aprile:2018dbl} and PandaX \\cite{Cui:2017nnn}. $\\mathbb{Z}_3$-symmetric models that stabilise the SM vacuum and produce the entire observed DM relic density are ruled out at 99\\% confidence; those constituting only a fraction of DM are ruled out at 98\\% confidence. The same is expected of other $\\mathbb{Z}_N$-symmetric models with $N>3$, which also feature non-self-adjoint DM.\n\nAs in the $\\mathbb{Z}_2$-symmetric case, a Bayesian analysis prefers the higher-mass part of the parameter space, due to the fine-tuning needed to achieve agreement with all experimental data in both the resonance and semi-annihilation (intermediate mass) regions. In this case, the additional tuning in $\\mu_3$ required to satisfy the condition $\\mu_3 \\leq 2\\sqrt{\\lambda_{\\sss S}}m_{\\sss S}$ -- and to achieve sufficient semi-annihilation in the intermediate-mass region -- further penalises these regions.\n\n\\begin{figure}[t]\n\t\\centering\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_VDM_2_1_like2D_combined_low_mass}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_VDM_2_1_like2D_combined_high_mass}\n \\caption{Profile likelihoods of parameters in the $\\mathbb{Z}_2$-symmetric vector singlet Higgs portal dark matter model, including constraints from direct and indirect detection, the relic density of dark matter and LHC searches for invisible decays of the Higgs boson, along with various Standard Model, dark matter halo and nuclear uncertainties. \\textit{Left}: the low-mass resonance region. \\textit{Right}: the full mass range. Grey shading indicates the area that fails the unitarity cut (Eq.\\ \\ref{eq:vec_unitarity}). Orange annotations indicate the edge of the allowed parameter space along which the model reproduces the entire cosmological abundance of dark matter. Contours show 1 and 2$\\sigma$ confidence regions. White stars indicate the location of the best-fit point. From \\protect\\cite{HP}.}\n\t\\label{fig:vector}\n\\end{figure}\n\n\n\\subsubsection{$\\mathbb{Z}_2$-symmetric vector singlet}\n\nIf DM is a $\\mathbb{Z}_2$-symmetric vector singlet $V_\\mu$ interacting with the SM via the Higgs portal, its effective Lagrangian takes the form\n\\begin{equation}\n\\mathcal{L_V} = -\\frac{1}{4} W_{\\mu\\nu} W^{\\mu\\nu} + \\frac{1}{2} \\mu_V^2 V_\\mu V^\\mu - \\frac{1}{4!} \\lambda_{V} (V_\\mu V^\\mu)^2 + \\frac{1}{2} \\lambda_{hV} V_\\mu V^\\mu H^\\dagger H.\n\\label{eq:Lag_V}\n\\end{equation}\nHere $W_{\\mu\\nu} \\equiv \\partial_\\mu V_\\nu - \\partial_\\nu V_\\mu$ is the field strength tensor for the new vector. The tree-level DM mass has exactly the same form as Eq.\\ \\ref{ms}. Although all terms here are dimension 4, the theory is not renormalisable, as it possesses an explicit mass term for $V_\\mu$. Perturbative unitarity is violated at energies above this mass. In the \\textsf{GAMBIT}\\xspace analysis \\cite{HP}, this issue was avoided by excluding the region of parameter space\n\\begin{equation}\n0 \\le \\lambda_{hV} \\le 2m_V^2\/v_0^2\n\\label{eq:vec_unitarity}\n\\end{equation}\nfrom the analysis.\n\nThe phenomenology of the vector model is very similar to that of the $\\mathbb{Z}_2$-symmetric scalar variant, with the only major difference being the absence of the intermediate-mass solution due to the unitarity requirement (Fig.\\ \\ref{fig:vector}; the region excluded from the analysis due to the unitarity condition is shown in grey). The Bayesian analysis once again prefers the high-mass region due to the fine-tuning of nuisance parameters required in the resonance region.\n\n\\subsubsection{$\\mathbb{Z}_2$-symmetric Dirac \\& Majorana fermionic singlets}\n\nThe Lagrangians of the fermionic singlet Higgs portal models are\n\\begin{align}\n \\mathcal{L}_{\\chi} &= \\frac{1}{2} \\overline{\\chi} (i\\slashed{\\partial} - \\mu_\\chi) \\chi - \\frac{1}{2}\\frac{\\lambda_{h\\chi}}{\\Lambda_\\chi} \\Big(\\cos\\theta \\, \\overline{\\chi}\\chi + \\sin\\theta \\, \\overline{\\chi}i\\gamma_5 \\chi \\Big) H^\\dagger H,\n \\label{eq:Lag_chi}\\\\\n \\mathcal{L}_{\\psi} &= \\overline{\\psi} (i \\slashed{\\partial} - \\mu_\\psi) \\psi \\nonumber - \\frac{\\lambda_{h\\psi}}{\\Lambda_\\psi} \\Big(\\cos\\theta \\, \\overline{\\psi}\\psi + \\sin\\theta \\, \\overline{\\psi}i\\gamma_5 \\psi \\Big) H^\\dagger H,\n \\label{eq:Lag_psi}\n\\end{align}\nwith the Majorana variant denoted $\\chi$ and the Dirac variant $\\psi$. These noticeably possess dimension-5 effective portal operators suppressed by the scale of new physics $\\Lambda$, with both scalar ($CP$-even) and pseudoscalar ($CP$-odd) couplings. The degree to which the portal interaction violates $CP$ is dictated by the mixing angle $\\theta$, where $\\theta = 0$ corresponds to pure $CP$ conservation and $\\theta = \\frac\\pi2$ to maximal $CP$ violation.\n\nAs in the scalar and vector models, the portal interaction produces terms quadratic in the DM field following electroweak symmetry breaking. The pseudoscalar coupling leads to an imaginary mass term, which must be rotated away with the field transformation $X \\rightarrow e^{i\\gamma_5 \\alpha\/2} X$ for $X \\in \\{\\chi, \\psi\\}$, in order to arrive at the physical (real) mass. This introduces a new parameter $\\alpha$. The physical masses are then\n\\begin{equation}\nm_{X}^2 = \\left(\\mu_{X} + \\frac{1}{2}\\frac{\\lambda_{hX}}{\\Lambda_{X}} v_0^2 \\cos\\theta \\right)^2 + \\left(\\frac{1}{2}\\frac{\\lambda_{X}}{\\Lambda_{X}}v_0^2 \\sin\\theta \\right)^2.\n\\end{equation}\nThe rotation parameter $\\alpha$ is fixed by the requirement that the mass be real, so all phenomenology can be described by three parameters: $m_X$, $\\lambda_X\/\\Lambda_X$ and $\\xi \\equiv \\theta + \\alpha$. Notably, the pure $CP$-conserving theory ($\\theta = 0$) remains $CP$-conserving after electroweak symmetry breaking ($\\xi = 0$), but maximal $CP$ violation before electroweak symmetry breaking does not correspond to maximal violation after the symmetry is broken (i.e. $\\theta = \\frac\\pi2 \\notimplies \\xi = \\frac\\pi2$).\n\nWhilst the $CP$-even Higgs portal coupling leads to the familiar velocity and momentum-independent nuclear scattering cross-section, the $CP$-odd coupling gives rise to an interaction suppressed by $q^2$, the square of the momentum exchanged in the scattering event. This leads to an overall suppression of direct detection signals and corresponding constraints for $\\xi \\rightarrow \\frac\\pi2$. Conversely, the $CP$-odd coupling produces a velocity and momentum-independent annihilation cross-section, whereas the $CP$-even coupling gives rise to a velocity-suppressed annihilation cross-section.\n\n\\begin{figure}[tbp]\n\t\\centering\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_MDM_2_1_like2D_combined_low_mass}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_MDM_2_1_like2D_combined_high_mass}\\\\\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_MDM_2_3_like2D_combined}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_MDM_3_1_like2D_combined}\n \\caption{Profile likelihoods of parameters in the $\\mathbb{Z}_2$-symmetric Majorana fermion singlet Higgs portal dark matter model, including constraints from direct and indirect detection, the relic density of dark matter and LHC searches for invisible decays of the Higgs boson, along with various Standard Model, dark matter halo and nuclear uncertainties. The upper-left panel shows a zoomed-in view of the low-mass resonance region. Grey shading indicates the area that fails the unitarity cut (Eq.\\ \\ref{eq:fermion_unitarity}). Orange annotations indicate the edge of the allowed parameter space along which the model reproduces the entire cosmological abundance of dark matter. Contours show 1 and 2$\\sigma$ confidence regions. White stars indicate the location of the best-fit point. From \\protect\\cite{HP}.}\n\t\\label{fig:fermion1}\n\\end{figure}\n\n\\begin{figure}[tbp]\n\t\\centering\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_MDM_2_1_post2D_xi_free_TWalk_low_mass}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_MDM_2_1_post2D_xi_free_TWalk_high_mass}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_MDM_2_3_post2D_xi_free_TWalk}\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_MDM_3_1_post2D_xi_free_TWalk}\n \\caption{Posterior probability densities from a Bayesian analysis of the $\\mathbb{Z}_2$-symmetric Majorana fermion singlet Higgs portal dark matter model, using the same likelihood functions as Fig.\\ \\protect\\ref{fig:fermion1}. White bullets indicate posterior means; other annotations are as in Fig.\\ \\protect\\ref{fig:fermion1}. From \\protect\\cite{HP}.}\n\t\\label{fig:fermion2}\n\\end{figure}\n\nProfile likelihoods from the global fit to the Majorana fermion model are shown in Fig.\\ \\ref{fig:fermion1}. Results for Dirac fermion dark matter are broadly very similar, and differ from the Majorana case only in the exact location of the border of the allowed parameter space, reflecting the essentially inconsequential nature of the relative factors of 2 between the two Lagrangians. Grey regions correspond to the regime\n\\begin{equation}\n\\lambda_{hX}\/\\Lambda_{X} \\geq 2\\pi\/m_X,\n\\label{eq:fermion_unitarity}\n\\end{equation}\nwhere the validity of the EFT becomes questionable. Further discussion on this issue can be found in Ref.\\ \\cite{HP}; it would also be possible to unitarise the theory, and draw further constraints in this region, using the $K$-matrix formalism \\cite{Bell:2016obu,Balaji:2018qyo}.\n\nThe preferred regions in the mass-coupling plane (upper panels of Fig.\\ \\ref{fig:fermion1}) include the now-familiar resonance and high-mass regions. However, unlike the vector and scalar models, these are fully connected by valid models at all masses, with the preferred region bounded from below mostly by the relic density constraint, supported by indirect detection. This is because profiling over $\\xi$ allows for the selection of $CP$-violating couplings in order to avoid constraints from direct detection. The degree of tuning in $\\xi$ required to achieve this is apparent in the lower panels of Fig.\\ \\ref{fig:fermion1}, where it is clear that good fits can be found for any value of $\\xi$ in the resonance region, but that higher masses require some degree of $CP$ violation in order to avoid direct detection. This becomes even clearer in the equivalent Bayesian results shown in Fig.\\ \\ref{fig:fermion2}, where intermediate masses and couplings are disfavoured relative to other regions, due to the need to make $CP$ violation nearly maximal in order to avoid direct detection.\n\n\\begin{figure}[t]\n\t\\centering\n \\includegraphics[width=0.495\\columnwidth]{figures\/plot_MDM_3_post1D_xi_free_TWalk}\n \\caption{Marginalised one-dimensional posterior probability density for the $CP$-mixing parameter $\\xi$ in the $\\mathbb{Z}_2$-symmetric Majorana fermion singlet Higgs portal dark matter model. This result has been extracted from the same analysis as that shown in Fig.\\ \\protect\\ref{fig:fermion2}. The value $\\xi = 0 = \\pi$ corresponds to $CP$ conservation; a clear preference for violation of $CP$ symmetry is evident. The blue bullet indicates the posterior mean value of $\\xi$, and the red star the value of $\\xi$ at the best-fit sample. From \\cite{HP}.}\n\t\\label{fig:fermion3}\n\\end{figure}\n\nIntegrating the posterior over all parameters other than $\\xi$ (Fig.\\ \\ref{fig:fermion3}), there is a clear preference for $CP$ violation. This reflects the fact that the more $CP$ violation permitted, the broader the range of other parameters able to give good fits to the combined data of all experiments. Performing Bayesian model comparison between the full model and its pure $CP$-conserving subspace (i.e.\\ $\\xi = 0$) results in Bayes factors of between 70:1 and 140:1, depending on the adopted priors. This indicates a strong preference for $CP$ violation in fermionic Higgs portal models. Bayesian model comparison between the scalar, vector and fermionic variants of the Higgs portal DM model reveals essentially equal odds for each of the scalar and fermionic models, but a 6:1 preference for all of these models over the vector variant.\n\n\\subsection{Axions}\n\n\\subsubsection{Axion models and their implementation in \\textsf{GAMBIT}\\xspace}\n\nAxions are an intriguing theoretical possibility due to their ability to solve the strong-$CP$ problem of the SM whilst providing a credible DM candidate~\\cite{Preskill:1982cy,Abbott:1982af,Dine:1982ah,1986_turner_axiondensity}. One can also use axion-like particles to reconcile various tensions between astrophysical observations and theory, including the cooling of white dwarfs~\\cite{Isern:1992gia,1205.6180,1211.3389,1512.08108,1605.06458,1605.07668,1708.02111}, and the transparency of the Universe to gamma rays~\\cite{0707.4312,0712.2825,1001.0972,1106.1132,1201.4711,1302.1208}.\n\nThe strong-$CP$ problem is ultimately a fine-tuning problem, arising from the fact that the SM symmetries permit a $CP$-odd term in the SM Lagrangian density of the form:\n\n\\begin{equation}\n\\mrm{\\pazocal{L}}{QCD} \\supset - \\frac{\\mrm{\\alpha}{S}}{8\\pi} \\mrm{\\theta}{QCD} G_{\\mu \\nu}^a \\widetilde{G}^{\\mu \\nu, a} \\, , \\label{eq:QCDLagrangian}\n\\end{equation}\nwhere $G_{\\mu \\nu}^a$ is the gluon field strength tensor, $\\widetilde{G}^{\\mu \\nu, a}$ is its dual (both of which have the $SU(3)$ gauge index $a$ explicitly shown), and $\\mrm{\\alpha}{S}$ is the strong coupling constant. The angle $\\mrm{\\theta}{QCD} \\in [-\\pi, \\pi]$ is a free parameter. In the SM, the term also receives a contribution from the chiral anomaly which, for down- and up-type Yukawa matrices $Y_d$ and $Y_u$, replaces $\\mrm{\\theta}{QCD}$ by the effective angle\n\\begin{equation}\n\\mrm{\\theta}{eff} \\equiv \\mrm{\\theta}{QCD} - \\arg \\left [ \\det (Y_dY_u) \\right ] \\, . \\label{eq:thetaeff}\n\\end{equation}\nA non-zero $\\mrm{\\theta}{eff}$ would result in $CP$-violating effects in strong interactions, which are severely constrained by observed upper limits on the electric dipole moment of the neutron, demanding $|\\mrm{\\theta}{eff}|\\mathrel{\\rlap{\\lower4pt\\hbox{$\\sim$}}\\raise1pt\\hbox{$<$}} \\num{e-10}$ \\cite{1509.04411}. Naively, this can only be avoided in the SM by fine-tuning the value of $\\mrm{\\theta}{QCD}$ to cancel the contribution from the chiral anomaly.\n\nAn alternative solution, first proposed by Peccei and Quinn~\\cite{1977_pq_axion1,1977_pq_axion2}, is to add a new global, axial $U(1)$ symmetry spontaneously broken by the vacuum expectation value $v$ of a complex scalar field. This breaking has an associated pseudoscalar Nambu-Goldstone boson, $a(x)$, which supplements $\\mrm{\\theta}{eff}$ by a new term $Na(x)\/v$, where the non-zero integer $N$ is the colour anomaly of the added symmetry. The Vafa-Witten theorem~\\cite{1984_vafa_vafawitten1,1984_vafa_vafawitten2} can then be used to show that $\\mrm{\\theta}{eff}+Na(x)\/v$ is dynamically driven to zero, solving the strong $CP$ problem.\n\nIn the resulting theory of the QCD axion, the axion is practically massless until the time of the QCD phase transition, due to a shift symmetry of the $U(1)$ phase, which prevents a mass term in the Lagrangian. After this, however, it picks up a small, temperature-dependent mass due to breaking of the continuous shift symmetry by fluctuations of the gluon fields. This gives rise to an effective axion potential\n\\begin{equation}\n\tV(a) = f_a^2 \\, m_a^2 \\, \\left[1 - \\cos (a\/f_a) \\right] \\, , \\label{eq:axion_eff_pot}\n\\end{equation}\nwhere $m_a$ is the temperature-dependent axion mass and $f_a \\equiv v\/N$. The zero-temperature axion mass, $m_{a,0}$, can be calculated using next-to-leading order chiral perturbation theory, and it turns out to be inversely proportional to $f_a$ for the QCD axion. At higher temperatures, numerical estimates of the mass are available from lattice QCD results, which can be described to a good approximation by\n\\begin{equation}\nm_a(T) = m_{a,0}\n\\begin{cases}\n\\hfil 1 \\hfil & \\mathrm{if \\; } T \\leq T_\\chi \\\\\n\\left ( \\frac{T_\\chi}{T} \\right )^{\\beta\/2} & \\mathrm{otherwise}\n\\end{cases} \\, . \\label{eq:axionmass}\n\\end{equation}\n$T_\\chi$ and $\\beta$ are in principle calculable, but can be left as nuisance parameters in order to account for systematic uncertainties in the calculations.\n\nIn fact, QCD axions are only one instance of a general class of \\emph{axion-like particles} (ALPs), which could generally result from the breaking of a $U(1)$ symmetry at some scale $f_a$, with mass generation occurring via the explicit breaking of the residual symmetry at some lower scale $\\Lambda$~\\cite{1987_kim_lightpseudoscalars,1002.0329,1801.08127}. It can be shown that in a Friedmann-Robertson-Walker-{Lema\\^itre} universe, a QCD axion or ALP field $\\theta(t)=a(t)\/f_a$ satisfies the equation of motion\n\\begin{equation}\n \\ddot{\\theta} + 3H(t) \\, \\dot{\\theta} + m_a^2(t) \\, \\sin (\\theta) = 0,\n\\label{eq:AxionFieldEq}\n\\end{equation}\nwhere we have assumed the canonical axion potential of\n\\begin{equation}\n V(\\theta) = f_a^2m_a^2\\left[1-\\cos (\\theta)\\right]. \\label{eq:potential}\n\\end{equation}\nThis is subject to the boundary condition $\\theta(\\mrm{t}{i}) = \\mrm{\\theta}{i}$ and $\\dot{\\theta}(\\mrm{t}{i}) = 0$, where $\\mrm{\\theta}{i}$ is called the \\textit{initial misalignment angle}.\n\nThe \\textsf{GAMBIT}\\xspace collaboration completed a comprehensive study of axion and broader ALP theories in 2018~\\cite{Axions}, using an extensive list of experimental constraints. These rely on the interactions of ALPs with SM matter, which can be studied in an effective field theory framework~\\cite{Kaplan:1985dv,1985_srednicki_axioneft,1986_georgi_axioneft}.\n\nThe most general axion\/ALP model in \\textsf{GAMBIT}\\xspace assumes the effective Lagrangian density to take the form\n\\begin{equation}\n\t\\pazocal{L}_a^\\mathrm{int} = -\\frac{f_ag_{a\\gamma\\gamma}}{4} \\theta F_{\\mu\\nu}\\widetilde{F}^{\\mu\\nu} - \\frac{f_ag_{aee}}{2m_e} \\bar{e}\\gamma^\\mu\\gamma_5e\\partial_\\mu \\theta \\, .\\label{eq:ax:lagrange}\n\\end{equation}\nNote that this provides for possible axion-photon and axion-electron interactions, whilst ignoring terms for other interactions that do not currently give rise to interesting experimental observables. The complete family tree of \\textsf{GAMBIT}\\xspace axion\/ALP models is shown in Fig.\\ \\ref{fig:AxionModelTree}, headed by the \\textsf{GeneralALP}\\xspace model, whose parameters have all now been defined. This provides a phenomenological description of axion physics that is not constrained to give physical solutions, as the couplings are not inversely proportional to $f_a$.\n\nThe \\textsf{QCDAxion}\\xspace model appears as a child model, and differs from the more general case by having tight constraints on some parameters, arising from the known relationships with the QCD scale. The axion-electron coupling is traded for the model-dependent form factor~$C_{aee}$\n\\begin{equation}\n\tg_{aee} = \\frac{m_e}{f_a} \\; C_{aee}\\, , \\label{eq:qcdaxioncouplings1}\n\\end{equation}\nwhilst the axion-photon coupling is replaced by the model-dependent ratio of the electromagnetic and colour anomalies~$E\/N$\n\\begin{equation}\n\tg_{a\\gamma\\gamma} = \\frac{\\mrm{\\alpha}{EM}}{2\\pi f_a}\\left(\\frac{E}{N} - \\widetilde{C}_{a\\gamma\\gamma}\\right) \\, .\\label{eq:qcdaxioncouplings2}\n\\end{equation}\n$\\widetilde{C}_{a\\gamma\\gamma}$ is a model-independent contribution from axion-pion mixing, which is taken from Ref.~\\cite{1511.02867}, and assigned a nuisance likelihood with a relevant uncertainty. Note that the ratio $E\/N$ should in principle take discrete values, but it is sampled as a continuous parameter for convenience, seeing as the possible rational values that it can take are close together. The final nuisance parameter of the \\textsf{QCDAxion}\\xspace model is $\\Lambda_\\chi$, which results from replacing the parameter $m_{a,0}$ of the \\textsf{GeneralALP}\\xspace model by an energy scale such that\n\\begin{equation}\n\tm_{a,0} \\equiv \\frac{\\Lambda_\\chi^2}{f_a} \\, .\n\\end{equation}\nThe value of $\\Lambda_\\chi$ is taken from first-principle calculations of the zero-temperature axion mass provided in Ref.~\\cite{1511.02867},\\footnote{This value was later updated in \\cite{Gorghetto:2018ocs}, after the appearance of Ref.\\ \\cite{Axions}.} and it is subject to a Gaussian nuisance likelihood.\n\nThe other models of interest for this review are the \\textsf{KSVZAxion}\\xspace and \\textsf{DFSZAxion}\\xspace model variants, which involve further field content being added to the SM. In \\textsf{KSVZAxion}\\xspace models~\\cite{1979_kim_ksvz,1980_shifman_ksvz}, the SM is supplemented by one or more electrically neutral, heavy quarks, and there are no tree-level interactions between the axion and SM fermions. There is, however, still an axion-photon interaction, which generates an axion-electron interaction at one loop. The \\textsf{GAMBIT}\\xspace study investigated four different \\textsf{KSVZAxion}\\xspace models, distinguished only by the choice of $E\/N$ from the set 0, 2\/3, 5\/3 and 8\/3.\n\n\\textsf{DFSZAxion}\\xspace models supplement the SM by an additional Higgs doublet~\\cite{1980_zhitnitsky_dfsz,1981_dine_dfsz}, which results in direct axion-electron interactions. Defining the ratio of the two Higgs vacuum expectation values to be $\\tan (\\beta^\\prime)$, one can write two variants of the \\textsf{DFSZAxion}\\xspace scenario as\n\\begin{align}\n\t\\begin{array}{lll}\n\t\tC_{aee} = \\sin^2 (\\beta^\\prime)\\left\/3\\right., \\quad \\phantom{.} & E\/N = 8\/3 \\quad \\phantom{.}& (\\textsf{DFSZAxion-I}\\xspace) \\, \\\\\n\t\tC_{aee} = \\left[1-\\sin^2 (\\beta^\\prime) \\right]\\left\/3\\right., \\quad \\phantom{.} & E\/N = 2\/3 \\quad \\phantom{.} & (\\textsf{DFSZAxion-II}\\xspace) \\,\n\t\\end{array} \\label{eq:dfsz:caee}.\n\\end{align}\nIt is thus convenient to replace the parameter~$C_{aee}$ in the \\textsf{QCDAxion}\\xspace model by $\\tan (\\beta^\\prime)$.\n\n\\begin{figure}[bt]\n\t\\centering\n\t\\input{include\/tree.tex}\n\t\\caption{Family tree of axion models in \\textsf{GAMBIT}\\xspace. The numbers in brackets refer to the number of model parameters; $(n+m)$ indicates $n$ (largely unconstrained) fundamental parameters of the model and $m$ (typically well-constrained) nuisance parameters. From \\cite{Axions}.}\n\t\\label{fig:AxionModelTree}\n\\end{figure}\n\n\\subsubsection{Experimental constraints on axions}\n\\label{sec:axionL}\nMany experiments are sensitive to the axion theories described here, and current null results place tight constraints on axions for specific combinations of masses and coupling strengths. Here we provide a brief review of those constraints, referring the reader to Ref.~\\cite{Axions} for a detailed description of the experimental likelihoods.\n\n\\begin{itemize}\n\\item \\textbf{Light-shining-through-wall (LSW) experiments: }Photon-axion interactions would allow photons to pass through a wall by becoming an axion, only to convert back to a photon on the other side. LSW experiments attempt to observe this by shining laser light onto an opaque material in the presence of a strong magnetic field. The \\textsf{GAMBIT}\\xspace LSW likelihood uses the results from the ALPS-I experiment, using data for both evacuated and gas-filled magnets~\\cite{1004.1313}.\n\\item \\textbf{Helioscopes: }Axion production in the Sun can be probed by observing the solar disc with a long magnet contained in an opaque casing. Any axions produced in the Sun that made it to Earth would pass through the exterior, and potentially convert to photons within the magnetic field in the interior. The details of solar axion production depend on the solar model, in addition to the axion-photon and axion-electron couplings. The \\textsf{GAMBIT}\\xspace axion studies utilise the AGSS09met solar model~\\cite{Serenelli09,AGSS} and its more recent iteration~\\cite{1611.09867}, and utilise two separate likelihoods for the 2007 and 2017 results of the CAST experiment~\\cite{hep-ex\/0702006,1705.02290}.\n\\item \\textbf{Haloscopes (cavity experiments): }Axion haloscopes aim to detect resonant axion-photon conversion inside a tunable cavity~\\cite{1983_Sikivie, 1985_Sikivie}, with microwave cavities providing the greatest current sensitivity to axions. Unfortunately, the resonant nature of the experiment means that one obtains highly sensitive constraints only within a very narrow mass range. The ability of haloscope experiments to detect axions depends on their cosmological abundance, as well as the galactic DM velocity distribution~\\cite{2011_Hoskins}. The \\textsf{GAMBIT}\\xspace study combines separate likelihood terms for the Rochester-Brookhaven-Fermi (RBF)~\\cite{DePanfilis:1987dk,Wuensch:1989sa}, University of Florida (UF)~\\cite{Hagmann:1990tj}, ADMX 1998-2009~\\cite{astro-ph\/9801286,Asztalos:2001tf,astro-ph\/0310042,astro-ph\/0603108,0910.5914} and ADMX 2018~\\cite{1804.05750} datasets.\n\\item \\textbf{Dark matter relic density: }Although axions are not a thermal relic such as those encountered in WIMP models, the relic abundance of axion DM is calculable numerically via the details of the realignment mechanism that follow from the equation of motion given in Eq~\\ref{eq:AxionFieldEq}. This can be compared with the observed value from the most recent \\emph{Planck} analysis~\\cite{Planck15cosmo}. The \\textsf{GAMBIT}\\xspace axion study applied this as both an upper limit (in which case axions are allowed to provide only a component of DM) and, in separate analyses of each model, a measurement. In the former case, anticipated yields in experiments that rely on the local DM density were scaled accordingly.\n\\item \\textbf{Distortions of gamma-ray spectra: }Axion-photon conversions could occur in strong galactic or inter-galactic magnetic fields, resulting in a distortion of the spectra of distant sources~\\cite{Raffelt:1987im,hep-ph\/0111311,hep-ph\/0204216,0704.3044}. There is a critical energy scale $E_{crit}$ at which photons will efficiently convert into axions, and it can be shown that spectral distortions only occur in real measurements when the critical energy lies within the spectral window of the instrument~\\cite{1205.6428,1305.2114}. This has the effect of localising constraints from spectral distortion measurements to specific ranges of the axion mass. The \\textsf{GAMBIT}\\xspace axion study utilises a likelihood based on H.E.S.S studies of the active galactic nucleus PKS 2155-304~\\cite{1311.3148}.\n\\item \\textbf{Supernova 1987A: }If axions had been produced in the SN1987A supernova explosion, they could have been converted to photons in the Galactic magnetic field, and detected as a coincident gamma ray burst by the Solar Maximum Mission~\\cite{Chupp:1989kx}. The absence of this observation has been used to constrain axion properties. The \\textsf{GAMBIT}\\xspace study uses a likelihood based on Ref.~\\cite{1410.3747}.\n\\item \\textbf{Horizontal Branch stars and the R parameter: }The existence of axions would provide an extra mechanism of energy loss for stars, causing them to cool faster~\\cite{Sato:1975vy,Raffelt:1990yz,book_raffelt_laboratories}. This would affect the relative time that stars spend on the Horizontal Branch (HB) and upper Red Giant Branch (RGB), which in turn sets the observed ratio of the numbers of stars on these branches ($R = \\mrm{N}{HB}\/\\mrm{N}{RGB}$). Theory suggests that axions would have the most significant impact on the lifetimes of HB stars, leading to a reduction in $R$. The \\textsf{GAMBIT}\\xspace $R$ parameter likelihood is based on the comparison of a calculation of the $R$ parameter for axion theories~\\cite{1512.08108,1983A&A...128...94B,Raffelt:1989xu,1311.1669,1406.6053} with the observed value of $\\mrm{R}{obs}=1.39 \\pm 0.03$~\\cite{1406.6053}, which is based on a weighted average of cluster count obervations \\cite{astro-ph\/0403600}.\n\\item \\textbf{White Dwarf cooling hints: }White dwarfs (WDs) are intriguing axion laboratories for several reasons. The first is that energy loss via axion production in WDs can be probed experimentally by using measurements of the oscillations of their radii and luminosities. These can be related to their internal structure via astroseismology, and measurements of the decrease in the oscillation periods can be related to energy loss. The second reason is that WDs have electron-degenerate cores, allowing us to probe the axion-electron coupling rather than the electron-photon coupling. A number of previous studies have calculated the expected period decrease in the presence of axions. The \\textsf{GAMBIT}\\xspace WD cooling likelihood is based on interpolation of the results and uncertainties found in Refs.~\\cite{1205.6180,1211.3389,1605.06458,1605.07668}. Current evidence suggests that WDs actually require an additional cooling mechanism relative to standard models, but this remains controversial due to a number of experimental and theoretical issues. The \\textsf{GAMBIT}\\xspace axion paper thus contains studies generated both with and without WD cooling hints added to the combined likelihood.\n\\end{itemize}\n\n\\subsubsection{\\textsf{GAMBIT}\\xspace results for the \\textsf{QCDAxion}\\xspace model}\n\nAlthough the \\textsf{GAMBIT}\\xspace axion paper contained results for all of the models described above, we will here concentrate on the \\textsf{QCDAxion}\\xspace results in the interests of brevity. The various parameters (including nuisance parameters) are shown in Table~\\ref{tab:priors:QCDAxion}, along with the chosen priors and prior ranges. For each of the nuisance parameters, the prior range is chosen to cover a range of approximately $-5\\sigma$ to $+5\\sigma$ around the known central value, where $\\sigma$ is the known uncertainty. The range of $E\/N$ values is selected to encompass those encountered in a broad range of previous axion model studies, whilst the range on $f_a$ is driven by the requirement that the range of possible axion masses reaches from very small masses to the the largest mass allowed by bounds on hot DM. Our choice of a log prior for $f_a$ is motivated by the fact that the scale is unknown. $C_{aee}$ is explored in a generous range around 1, whilst the causal structure of the early Universe motivates our use of a flat prior on the initial misalignment angle $\\mrm{\\theta}{i}$. The local DM density $\\rho_0$ is given the same treatment as in previous \\textsf{GAMBIT}\\xspace studies.\n\n\\begin{table}[tbp]\n\t\\caption{Prior choices for \\textsf{QCDAxion}\\xspace models in \\cite{Axions}.\\label{tab:priors:QCDAxion}}\n\t\\footnotesize\n\t\\centering\n\t\\begin{tabular}{@{}lccc}\n\t\t\\toprule\n\t\t\\textbf{Model} & \\multicolumn{2}{l}{\\textbf{Parameter range\/value}} & \\textbf{Prior type} \\\\\n\t\t\\midrule\n\t\t\\textsf{QCDAxion}\\xspace & \\iuo{f_a}{\\text{G\\eV}\\xspace} & \\prrange{e6}{e16} & log \\\\\n\t\t& \\iuo{\\Lambda_\\chi}{\\text{M\\eV}\\xspace}& \\prrange{73}{78} & flat \\\\\n\t\t& $\\widetilde{C}_{a\\gamma\\gamma}$ & \\prrange{1.72}{2.12} & flat \\\\\n\t\t& $E\/N$ & \\prrange{-1.33333}{174.667} & flat \\\\\n\t\t& $C_{aee}$ & \\prrange{e-4}{e4} & log \\\\\n\t\t& $\\mrm{\\theta}{i}$ & \\prrange{-3.14159}{3.14159} & flat \\\\\n\t\t& $\\beta$ & \\prrange{7.7}{8.2} & flat \\\\\n\t\t& \\iuo{T_\\chi}{\\text{M\\eV}\\xspace}& \\prrange{143}{151} & flat \\\\\n\t\t\\midrule\n\t\tLocal DM density & \\iuo{\\rho_0}{\\text{G\\eV}\\xspace\\per\\centi\\metre^3} & \\prrange{0.2}{0.8} & flat \\\\\n\t\t\\bottomrule\n\t\\end{tabular}\n\\end{table}\n\n\\begin{figure}[tbp]\n\t\\centering\n\t{\n\t\t\\includegraphics[width=0.49\\linewidth]{figures\/diver\/plot_2_10M1_QCDAxion_100_6_like2D}\n\t\t\\hfill\n\t\t\\includegraphics[width=0.49\\linewidth]{figures\/diver\/plot_2_10M2_QCDAxion_100_6_like2D}\n\t}\n\t{\n\t\t\\includegraphics[width=0.49\\linewidth]{figures\/diver\/plot_2_10M1_QCDAxion_100_108_like2D}\n\t\t\\hfill\n\t\t\\includegraphics[width=0.49\\linewidth]{figures\/diver\/plot_2_10M2_QCDAxion_100_108_like2D}\n\t}\n\t\\caption{Profile likelihoods~(from \\diver) for \\textsf{QCDAxion}\\xspace models with upper limits~(\\textit{left}) and matching condition~(\\textit{right}) for the observed DM relic density. The upper and lower panels show the constraints on the anomaly ratio, $E\/N$, and the absolute value of the initial misalignment angle, $|\\mrm{\\theta}{i}|$, respectively. From \\cite{Axions}. \\label{fig:QCDAxion:frequentist}}\n\\end{figure}\n\nFig.\\ \\ref{fig:QCDAxion:frequentist} shows profile likelihood distributions in various planes of the \\textsf{QCDAxion}\\xspace parameters, obtained without the presence of WD cooling hints in the combined likelihood. The left panels show the result of imposing the relic density constraint as an upper limit. The exclusion of the low-$f_a$ (high mass) region, except at very low values of the axion-photon coupling (which is related to $E\/N$), arises from the $R$ parameter and CAST results. The slight reduction in the profile likelihood for masses lower than approximately \\SI{0.1}{\\micro\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace} also comes from the $R$ parameter likelihood; for such masses, the maximum allowed value for the axion-electron coupling ($C_{aee} \\leq \\num{e4}$) is not large enough to perfectly satisfy the $R$-parameter constraint. If axions are assumed to saturate the relic abundance of DM (right top panel), the high-mass region is excluded entirely due to the fact that the realignment mechanism cannot produce enough DM. The bottom row of Fig.\\ \\ref{fig:QCDAxion:frequentist} shows the allowed values for the initial misalignment angle. In the case that axions supply all of DM, we recover the familiar result that $\\left|\\mrm{\\theta}{i}\\right| \\ll 1$, for QCD axion masses of $m_{a,0} \\mathrel{\\rlap{\\lower4pt\\hbox{$\\sim$}}\\raise1pt\\hbox{$<$}} \\SI{0.1}{\\micro\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace}$, a fine-tuning that we will discuss further in the context of a Bayesian analysis.\n\nIn the left panel of Fig.\\ \\ref{fig:cooling:QCDAxion:RelicDens:twoD}, we show the marginalised Bayesian posterior in the $\\Omega_a h^2-m_{a,0}$ plane without WD cooling hints, demonstrating that the scan can find viable parts of the parameter space where axions consistent with all current experimental observations can account for a sizable fraction of dark matter. The situation is similar even when WD cooling hints are included (not shown). One can also observe an interesting bound on the axion mass. If the DM relic density constraint is applied as an upper limit, we find $\\SI{0.73}{\\micro\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace} \\le m_{a,0} \\le \\SI{6.1}{\\milli\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace}$ at 95\\% credibility (equal-tailed interval). Meanwhile, if axions must provide all of the observed dark matter, this changes to $\\SI{0.53}{\\micro\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace} \\le m_{a,0} \\le \\SI{0.13}{\\milli\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace}$.\n\n\\begin{figure}\n \\centering\n {\n \\includegraphics[width=0.49\\linewidth]{figures\/twalk\/plot_2_3042_QCDAxion_100_10_post2D}\n \\includegraphics[width=0.49\\linewidth]{figures\/twalk\/plot_2_3141_QCDAxion_100_107_post2D}\n }\n \\caption{Marginalised posterior for \\textsf{QCDAxion}\\xspace models with the DM relic density constraint treated as an upper limit. \\textit{Left}: constraints on the energy density in axions today, $\\Omega_a h^2$, without the inclusion of WD cooling hints. \\textit{Right}: constraints on the absolute value of the axion-photon coupling, $|g_{a\\gamma\\gamma}|$. This panel also includes WD cooling hints, but they have little impact on the result. For comparison, the right panel also shows the region for which QCD axions are not theoretically possible (red line and shading), as well as the frequentist $2\\sigma$ C.L. constraints on more general ALP models (dashed lines). From \\cite{Axions}. \\label{fig:cooling:QCDAxion:RelicDens:twoD}}\n\\end{figure}\n\nIn the right panel of Fig.\\ \\ref{fig:cooling:QCDAxion:RelicDens:twoD}, we show the marginalised posterior in the $|g_{a\\gamma\\gamma}|-m_{a,0}$ plane with the DM relic density constraint applied as an upper limit, and with WD cooling hints included. Also shown are the {na\\\"ive} bounds on the parameter space that result from phenomenological constraints on \\textsf{GeneralALP}\\xspace models and the maximum value of $E\/N$. The shape of the preferred region is partly formed by the effect of fine-tuning. At low axion masses, this is required to avoid dark matter overproduction, whilst at large axion masses it is required to achieve low values of $|g_{a\\gamma\\gamma}|$ through cancellations between $E\/N$ and $\\tilde{C}_{a\\gamma\\gamma}$. The preferred parameter region is localised within a few orders of magnitude in mass around $m_{a,0} \\sim \\SI{100}{\\micro\\ensuremath{\\text{e}\\mspace{-0.8mu}\\text{V}}\\xspace}$ and $g_{a\\gamma\\gamma} \\sim \\SI{e-12}{\\text{G\\eV}\\xspace^{-1}}$.\n\n\\begin{table}[tbp]\n\t\\caption{Prior choices for \\textsf{DFSZAxion-I}\\xspace, \\textsf{DFSZAxion-II}\\xspace and \\textsf{KSVZAxion}\\xspace models in \\cite{Axions}. Note that the priors listed in the first section of the table apply to all three models.\\label{tab:priors:DFSZvsKSVZ}}\n\t\\footnotesize\n\t\\centering\n\t\\begin{tabular}{lcccl}\n\t\t\\toprule\n\t\t\\textbf{Model} & \\multicolumn{2}{l}{\\textbf{Parameter range\/value}} & \\textbf{Prior type} & \\textbf{Comments}\\\\\n\t\t\\midrule\n\t\t& \\iuo{f_a}{\\text{G\\eV}\\xspace} & \\prrange{e6}{e16} & log & Applies to all\\\\\n\t\t& \\iuo{\\Lambda_\\chi}{\\text{M\\eV}\\xspace}& \\prrange{73}{78} & flat & Applies to all\\\\\n\t\t& $\\widetilde{C}_{a\\gamma\\gamma}$ & \\prrange{1.72}{2.12} & flat & Applies to all\\\\\n\t\t& $\\mrm{\\theta}{i}$ & \\prrange{-3.14159}{3.14159} & flat & Applies to all\\\\\n\t\t& $\\beta$ & \\prrange{7.7}{8.2} & flat & Applies to all\\\\\n\t\t& \\iuo{T_\\chi}{\\text{M\\eV}\\xspace}& \\prrange{143}{151} & flat & Applies to all\\\\\n\t\t\\midrule\n\t\t\\textsf{DFSZAxion-I}\\xspace & $E\/N$ & $8\/3$ & delta & \\\\\n\t\t& $\\tan(\\beta')$ & \\prrange{0.28}{140.0} & log & \\\\\n\t\t\\textsf{DFSZAxion-II}\\xspace & $E\/N$ & $2\/3$ & delta & \\\\\n\t\t& $\\tan(\\beta^\\prime)$ & \\prrange{0.28}{140.0} & log & \\\\\n\t\t\\textsf{KSVZAxion}\\xspace & $E\/N$ & $0$, $2\/3$, $5\/3$, $8\/3$ & delta & Discrete\\\\\n\t\t\\midrule\n\t\tLocal DM density & \\iuo{\\rho_0}{\\text{G\\eV}\\xspace\\per\\centi\\metre^3} & \\prrange{0.2}{0.8} & flat \\\\\n\t\t\\bottomrule\n\t\\end{tabular}\n\\end{table}\n\nThe \\textsf{GAMBIT}\\xspace Bayesian analysis of axion models also includes a model comparison of the \\textsf{QCDAxion}\\xspace, \\textsf{DFSZAxion-I}\\xspace, \\textsf{DFSZAxion-II}\\xspace and \\textsf{KSVZAxion}\\xspace, based on scans of the latter models that use the priors defined in Table~\\ref{tab:priors:DFSZvsKSVZ}. Bayesian evidence values $\\mathbb{Z}(\\mathcal{M})$ for each model $\\mathcal{M}$ were calculated using the \\MultiNest nested sampling package, before constructing the Bayes factor \\cite{Jeffreys:1939xee,10.2307\/2291091,10.2307\/4356165}\n\\begin{equation}\n\t\\mathcal{B} \\equiv \\frac{\\mathbb{Z}(\\mathcal{M}_1)}{\\mathbb{Z}(\\mathcal{M}_2)} \\equiv \\frac{\\int \\! \\mathcal{L}\\left(\\text{data}\\left| \\right. \\boldsymbol{\\theta_1} \\right) \\pi_1(\\boldsymbol{\\theta_1}) \\, \\mathrm{d} \\boldsymbol{\\theta_1}}{\\int \\! \\mathcal{L}\\left(\\text{data}\\left| \\right. \\boldsymbol{\\theta_2} \\right) \\pi_2(\\boldsymbol{\\theta_1}) \\, \\mathrm{d} \\boldsymbol{\\theta_2}} \\, ,\n\\end{equation}\nwhich relates two models $\\mathcal{M}_1$ and $\\mathcal{M}_2$ with parameters $\\boldsymbol{\\theta_1}$ and~$\\boldsymbol{\\theta_2}$. $\\pi_1$ and $\\pi_2$ are the priors on the parameters of the two models, and $\\mathcal{L}$ is the likelihood. The Bayes factor is connected to the ratio of posterior probabilities of the models being correct\n\\begin{equation}\n\t\\frac{\\mathcal{P}\\left(\\mathcal{M}_1 \\left| \\text{data} \\right.\\right)}{\\mathcal{P}\\left(\\mathcal{M}_2 \\left| \\text{data} \\right.\\right)} = \\mathcal{B} \\; \\frac{\\pi(\\mathcal{M}_1)}{\\pi(\\mathcal{M}_2)},\n\\end{equation}\nwhere the prior probabilities of the models themselves being correct are given by $\\pi(\\mathcal{M}_1)$ and $\\pi(\\mathcal{M}_2)$. In the following, it is assumed that $\\pi(\\mathcal{M}_1)=\\pi(\\mathcal{M}_2)$, causing the the posterior odds ratio to be equal to the Bayes factor.\n\nWithout cooling hints, the odds ratios for pairs of models provide insufficient evidence to favour any particular scenario. However, if it is demanded that axions solve the DM and WD cooling problems simultaneously, there is a positive preference for the \\textsf{QCDAxion}\\xspace model over the DFSZ- and KSVZ-type models, at a level of about 5:1. This results from the larger $C_{aee}$ values allowed in the \\textsf{QCDAxion}\\xspace model, which peaks at $C_{aee}\\approx 100$ in the one-dimensional marginalised posterior. Such a large coupling may cause a problem for model building. A frequentist analysis of the same scenario allows both the \\textsf{DFSZAxion}\\xspace and \\textsf{KSVZAxion}\\xspace models to be rejected with respect to the \\textsf{QCDAxion}\\xspace model with better than 99\\% confidence; if DM is instead allowed to consist only partially of axions, only \\textsf{KSVZAxion}\\xspace models can be rejected in this way.\n\n\\subsection{Right-handed neutrinos}\n\\subsubsection{Model definition}\nThe addition of right-handed ``sterile'' neutrinos to the SM has been proposed to explain the existence of neutrino flavour oscillations, which imply a non-zero neutrino mass. They also serve an aesthetic theoretical purpose, as neutrinos are the only elementary fermions in the SM to not have both left- and right-handed incarnations. Moreover, sterile neutrinos are permitted to have both a Dirac mass term $\\bar{\\nu_L}M_D\\nu_R$ and a Majorana mass term $\\bar{\\nu_R}M_M\\nu_R^c$, and specific choices of the latter allow sterile neutrinos to solve cosmological problems such as the baryon asymmetry of the Universe~\\cite{Fukugita:1986hr,Akhmedov:1998qx,Asaka:2005pn}, and the DM problem~\\cite{Dodelson:1993je, Shi:1998km}.\n\nA convenient parameterisation of a right-handed neutrino sector is the Casas-Ibarra parametrisation, amended to include 1-loop corrections to the left-handed neutrino mass matrix~\\cite{Casas:2001sr,Lopez-Pavon:2015cga}. This involves writing a matrix that encodes the mixing among left-handed neutrinos (LHNs) and right-handed neutrinos (RHNs) as\n\\begin{align}\n\\Theta = iU_{\\nu}\\sqrt{m_{\\nu}^\\text{diag}}\\mathcal{R}\\sqrt{\\tilde{M}^\\text{diag}}^{-1},\n\\label{CItheta}\n\\end{align}\nwhere $U_{\\nu}$ is the PMNS matrix, $m_{\\nu}^\\text{diag}$ is a diagonalised, one-loop-corrected LHN mass matrix and $\\tilde{M}^\\text{diag}$ is the analogous RHN mass matrix. $\\mathcal{R}$ is a complex, orthogonal matrix written as the product\n\\begin{align}\n\\mathcal{R} = \\mathcal{R}^{23}\\mathcal{R}^{13}\\mathcal{R}^{12}\\;,\n\\label{Rorder}\n\\end{align}\nwhere the $\\mathcal{R}^{ij}$ can, in turn, be parameterised by complex angles $\\omega_{ij}$ with\n\\begin{align}\n\\mathcal{R}^{ij}_{ii} &= \\mathcal{R}^{ij}_{jj} = \\cos\\omega_{ij}, \\\\\n\\mathcal{R}^{ij}_{ij} &= -\\mathcal{R}^{ij}_{ji} = \\sin\\omega_{ij}, \\\\\n\\mathcal{R}^{ij}_{kk} &= 1; k \\neq i,j\\;.\n\\end{align}\n\nWorking in the flavour basis in which the Yukawa couplings of the charged leptons are diagonal by construction, the PMNS matrix $U_\\nu$ can be written as\n\\begin{align}\\label{UnuParameterisation}\n U_{\\nu} = V^{23}U_{\\delta}V^{13}U_{-\\delta}V^{12}\\mathrm{diag}(e^{i\\alpha_1\/2},e^{i\\alpha_2\/2},1)\\;,\n\\end{align}\nwhere $U_{\\pm\\delta} = \\mathrm{diag}(e^{{\\mp}i\\delta\/2},1,e^{{\\pm}i\\delta\/2})$\nand $V^{ij}$ is parameterised by the LHN mixing angles $\\theta_{ij}$. The non-zero elements of $V^{ij}$ are analogous to those of $\\mathcal{R}$. $\\alpha_1$, $\\alpha_2$ and $\\delta$ are $CP$-violating phases that are not excluded \\emph{a priori}.\n\n\\begin{table}[t]\n \\centering\n \\begin{tabular}{l}\n \\toprule\n \\textbf{Parameter} \\\\\n \\midrule\n Active neutrino parameters\\\\\n \\quad$\\theta_{12}$ [rad] \\\\\n \\quad$\\theta_{23}$ [rad] \\\\\n \\quad$\\theta_{13}$ [rad] \\\\\n \\quad$m_{\\nu_0}$ [eV] \\\\\n \\quad$\\Delta m^2_{21}$ $[10^{-5}\\,\\text{eV}^2]$ \\\\\n \\quad$\\Delta m^2_{3l}$ $[10^{-3}\\,\\text{eV}^2]$ \\\\\n \\quad$\\alpha_1$, $\\alpha_2$ [rad] \\\\\n Sterile neutrino parameters\\\\\n \\quad$\\delta$ [rad] \\\\\n \\quad Re $\\omega_{ij}$ [rad] \\\\\n \\quad Im $\\omega_{ij}$ \\\\\n \\quad$M_I$ [GeV] \\\\\n \\quad$R_{\\rm{order}}$ \\\\\n Nuisance parameters \\\\\n \\quad$m_h$ [GeV]\\\\\n \\bottomrule\n \\end{tabular}\n \\caption{The full list of scanned parameters for the \\textsf{GAMBIT}\\xspace right-handed neutrino study \\cite{RHN}.}\n \\label{tab:scanpars}\n\\end{table}\n\n\\begin{figure}[tbp]\n \\centering\n \\includegraphics[width=0.45\\linewidth]{figures\/NH_M_Ue_capped.pdf}\n \\includegraphics[width=0.45\\linewidth]{figures\/IH_M_Ue_capped.pdf}\\\\\n \\includegraphics[width=0.45\\linewidth]{figures\/NH_M_Umu_capped.pdf}\n \\includegraphics[width=0.45\\linewidth]{figures\/IH_M_Umu_capped.pdf}\\\\\n \\includegraphics[width=0.45\\linewidth]{figures\/NH_M_Utau_capped.pdf}\n \\includegraphics[width=0.45\\linewidth]{figures\/IH_M_Utau_capped.pdf}\n \\caption{Profile likelihoods of right-handed neutrino models in the $M_I$ vs $U_{eI}^2$ (top), $M_I$ vs $U_{\\mu I}^2$ (middle) and $M_I$ vs $U_{\\tau I}^2$ (bottom) planes. Results are shown for normal (left) and inverted (right) neutrino mass ordering. From \\cite{RHN}.}\n \\label{fig:M_Ue_capped}\n\\end{figure}\n\n\\begin{figure}[tbp]\n \\centering\n \\includegraphics[width=0.45\\linewidth]{figures\/NH_M_U_capped.pdf}\n \\includegraphics[width=0.45\\linewidth]{figures\/IH_M_U_capped.pdf}\n \\caption{Profile likelihood of right-handed neutrino models in the $M_I$ vs $U_I^2$ plane for normal (left) and inverted (right) neutrino mass ordering. From \\cite{RHN}.}\n \\label{fig:M_U_capped}\n\\end{figure}\n\nA comprehensive, frequentist \\textsf{GAMBIT}\\xspace study of this scenario has recently been completed \\cite{RHN}. The full list of parameters considered in the \\textsf{GAMBIT}\\xspace RHN study is given in Table~\\ref{tab:scanpars}. Separate scans were done for the cases of a normal and an inverted mass hierarchy, and the scanning strategy used a number of carefully targeted scans to ensure that the regions near the various experimental bounds were convergently sampled. The scans used the full range of likelihoods implemented in \\textsf{NeutrinoBit}\\xspace, documented in Section~\\ref{sec:neutrinobit}, in addition to the extra routines described in \\textsf{FlavBit}\\xspace, \\textsf{DecayBit}\\xspace and \\textsf{PrecisionBit}\\xspace. The main analysis used a capped likelihood, in which each point is forced to have a likelihood which is equal to or worse than the SM ($\\mathcal{L} = \\min[\\mathcal{L}_{\\rm{SM}}, \\mathcal{L}_{\\rm{RHN}}]$). This is due to a number of excesses in individual experimental observations that --- although combining to give a small overall significance --- would bias the presentation of exclusion limits on RHN parameters. In this review, we concentrate on the resulting limits on RHNs, and direct the reader to the original study~\\cite{RHN} for a detailed discussion of the excesses.\n\nThe 1-loop Casas-Ibarra parameterisation used in the GAMBIT analysis \\cite{RHN} is valid for seesaw scenarios where the active-sterile mixing, $|\\Theta|^2$, is small. In principle, additional $|\\Theta|^4$ corrections could be expected in low-scale seesaw scenarios, and could only be captured by an exact expansion (e.g.\\ Schechter-Valle \\cite{Schechter:1980gr}). Nevertheless, the loss of generality in the Casa-Ibarra approximation is outweighed by its numerical and computational benefits. This parameterisation allows one to explicitly choose the masses of both active and sterile neutrinos, and is automatically consistent with oscillation data, allowing oscillation parameters to be easily treated as Gaussian nuisances. Alternative parameterisations would constitute a different effective prior on both the sterile and active neutrino parameters. Because the GAMBIT analysis is based on profile likelihoods, which are by construction prior-independent, switching parameterisation would only have the impact of making sampling less efficient, rather than causing any physical or statistical effect.\n\n\\subsubsection{RHN results}\n\nIn Fig.\\ \\ref{fig:M_Ue_capped}, we show, as functions of the heavy neutrino masses $M_I$, the constraints on the couplings $U^2_{\\alpha I}$ to the active neutrino flavours $\\alpha = (e,\\mu,\\tau)$; in Fig.\\ \\ref{fig:M_U_capped} we also show the overall constraints on their combination $U^2_I=\\sum_\\alpha U^2_{\\alpha I}$. The index $I$ can refer to any of the heavy neutrino flavours $I=(1,2,3)$, as their labelling has no physical significance. The profile likelihood is mostly flat at low values of the couplings, but exhibits characteristic drop-offs at higher values that result from specific experimental observations. The most dominant constraint varies with the RHN mass.\n\nAbove the masses of the weak gauge bosons, direct searches at colliders are not relevant, and the leading constraints on the RHN properties come from electroweak precision observables, CKM measurements and searches for lepton flavour violation (LFV). The upper limits on the $\\tau$ couplings are much larger than on the $e$ and $\\mu$ couplings, due to the fact that the EWPO and LFV limits are stronger for the $e$ and $\\mu$ flavours.\n\nWhen $M_I$ is between the $D$ meson and $W$ boson masses, direct search experiments dominate, as RHNs are efficiently produced via the $s$-channel exchange of on-shell $W$ bosons. The DELPHI and CMS results compete to impose the strongest bound in this region.\n\nBelow the $D$ meson mass, the dominant constraints come from direct searches at beam-dump experiments, in particular CHARM and NuTeV (above the kaon mass), PS-191 and E949 (between the pion and the kaon mass), and pion decay experiments at even lower mass. In the case of the $\\tau$ couplings, the direct search constraints are much weaker, and the most significant constraint instead comes from DELPHI searches for long-lived particles.\n\nFor $M_I$ values below 0.3\\,GeV, the global constraints are stronger than the sum of the individual contributions, due to an interplay between the lower bound from BBN, the upper bounds from direct searches and the constraints on RHN mixing from neutrino oscillation data (which disfavour large hierarchies amongst the couplings to individual SM flavours). The BBN lifetime constraint does not have an observable effect on the individual couplings, but it does force their combination to be greater than a certain value (as seen in Fig.\\ \\ref{fig:M_U_capped}).\n\nFinally, we point out that although this analysis included the active neutrino oscillation likelihoods contained in \\textsf{NeutrinoBit}, based on the results of \\textsf{NuFit} \\cite{NuFit15}, it would not change the results even slightly if one were to replace these with nuisance likelihoods from either of the other main 3-flavour neutrino fitting groups \\cite{deSalas:2017kay,Capozzi:2018ubv}. This is because the results of all three groups are highly consistent, and the preferred parameter region is where the approximate $B-L$ symmetry holds and the oscillation constraints are essentially irrelevant. As the fits allow $m_{\\nu_0}\\to 0$, there is no lower limit implied on $M_I$ from oscillation data, but rather only from BBN (the effects of which were modelled under the massless neutrino approximation).\n\n\n\\section{Summary}\n\\label{sec:summary}\n\n\\textsf{GAMBIT}\\xspace is an open-source software framework for combining all relevant experimental constraints on theories for physics Beyond the Standard Model. It includes extensive libraries of theory and likelihood routines for dark matter, flavour, collider, neutrino and precision observables, along with spectrum and decay calculations, all for a range of popular and highly plausible theories of new physics. It features an array of different statistical samplers, a hierarchical model database, an automated engine for building calculations based on graph theory, and the ability to connect to external physics calculators with ease.\n\nIn the two years since its release, \\textsf{GAMBIT}\\xspace has produced seven global analyses of leading theories for physics beyond the Standard Model. In supersymmetry, this includes analyses of the CMSSM, NUHM1, NUHM2, a 7-parameter weak-scale MSSM, and an electroweakino effective field theory known as the EWMSSSM. Results indicate a $3.3\\sigma$ combined preference in LHC searches for weak production of light charginos and neutralinos. In Higgs portal dark matter, \\textsf{GAMBIT}\\xspace results cover $\\mathbb{Z}_2$ and $\\mathbb{Z}_3$-symmetric scalar singlet models, as well as $\\mathbb{Z}_2$-symmetric vector and fermion models. All can provide good fits to experimental constraints, but fermionic models strongly prefer to violate $CP$. Scalar models can not only solve the dark matter problem, but can also stabilise the vacuum of the standard model if and only if they possess TeV-scale masses and respect a $\\mathbb{Z}_2$ symmetry. GAMBIT studies indicate that QCD axions are most likely to constitute a fraction of dark matter rather than the entire amount, and to reside in a mass window between about $10^{-7}$ and $10^{-3}$\\,eV. Right-handed neutrinos are constrained by a wide array of different searches at different masses; interactions with electrons and muons are the most strongly constrained, with constraints on couplings to tau leptons somewhat weaker.\n\n\\textsf{GAMBIT}\\xspace is a powerful tool for testing theories of new physics. The code can be obtained from \\url{https:\/\/gambit.hepforge.org}. All samples, input files and benchmark points resulting from the \\textsf{GAMBIT}\\xspace physics studies discussed in this review can also be obtained from \\textsf{Zenodo}, by visiting \\url{https:\/\/zenodo.org\/communities\/gambit-official}.\n\n\\section*{Acknowledgements}\nWe thank our collaborators within the \\textsf{GAMBIT}\\xspace Community for their essential and extensive contributions to the work reviewed in this article, and Tom\\'as Gonzalo in particular for comments on right-handed neutrinos. The majority of plots presented in the review were produced with \\pippi \\cite{pippi}. We acknowledge PRACE for awarding us access to Marconi at CINECA, Italy.\n\n\\bibliographystyle{JHEP_pat}\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{INTRODUCTION}\nVisual place recognition in changing environments is the problem of finding matchings between two sets of observations, database and query, despite severe appearance changes.\nIt is required for loop closure detection in SLAM (Simultaneous Localization and Mapping) and for candidate selection in image-based 6D localization systems.\nThe common pipeline for visual place recognition in changing environments is shown in Fig.~\\ref{fig:intro} (top):\nGiven the two sets of images, database (the ``map'') and query (later or current run), a descriptor is computed for each image.\nSubsequently, each descriptor of the database is compared pairwise to each descriptor of the query to get a similarity matrix $\\hat{S}^{DB\\times Q}$ that can finally be used to determine potential place matches.\nAs illustrated in Fig.~\\ref{fig:intro} (bottom), there are various additional approaches to improve the place recognition pipeline either by preprocessing the descriptors, e.g., with feature standardization, or by postprocessing the similarity matrix, e.g., with sequence search in the similarity matrix.\nIn addition to these pre- and postprocessing steps, additional information can be exploited like database image poses or the so far rarely used intra-database and intra-query similarities.\nHowever, there are only few methods that exploit such additional knowledge in order to further improve the place recognition performance.\n\nIn this paper, we address the problem of systematically exploiting additional information by proposing a versatile, expandable framework for different kinds of prior knowledge from or about database and query.\nSpecifically, we\n\\begin{itemize}\n \\item propose a versatile, expandable graph-based framework that formulates place recognition as a sparse non-linear least squares optimization problem in a factor graph (Sec.~\\ref{sec:algo})\n \\item discuss different sources of information in place recognition problems and propose a loop-rule and an exclusion-rule to exploit inherent structural properties of the place recognition problem (Sec.~\\ref{sec:binary_factor})\n \\item demonstrate how this framework can be used to integrate the different sources of information, e.g., we provide implementations of the loop- and exclusion-rules in terms of cost functions for factors. These either exploit intra-set descriptor similarities within database and\/or query, or, if available, additional knowledge about database image poses (Sec.~\\ref{sec:binary_factor})\n \\item demonstrate the versatility of the framework by proposing an n-ary factor that mimics the sequence processing approach of SeqSLAM \\cite{seqSLAM} (Sec.~\\ref{sec:nary_factor})\n \\item present the optimization using standard non-linear least squares optimization tools (Sec.~\\ref{sec:optim})\n \\item provide comprehensive experimental evaluation on a variety of datasets, configurations and existing methods from the literature (Sec.~\\ref{sec:results})\n\\end{itemize}\n\nThe paper closes with a discussion of the current implementation and potential extensions in Sec.~\\ref{sec:discussion}.\n\n\\begin{figure}[tb]\n \\centering \n \\includegraphics[width=1\\linewidth]{images\/pr_pipeline.pdf}\n \\vspace{-0.7cm}\n \\caption{The basic place recognition pipeline (above the horizontal dashed line) can be extended with additional information (below this line). Established approaches are standardization of descriptors and sequence processing. We propose a graph-based framework to integrate various sources of additional information in a common optimization formulation. An example for so far rarely used information are intra-set similarities.} \n \\vspace{-0.4cm}\n \\label{fig:intro}\n\\end{figure}\n\n\n\\section{RELATED WORK}\n\\label{sec:related_work}\nVisual place recognition in changing environments is a subject of active research.\nAn overview of existing methods is given in \\cite{Lowry2016}.\nIn the present paper, the basic source of information to match places are image descriptors.\nThe authors of \\cite{Suenderhauf2015} demonstrated that intermediate CNN-layer outputs like the \\textit{conv3}-layer from AlexNet \\cite{alexnet} trained for image classification can serve as holistic image descriptors to match places across condition changes between database and query.\nMoreover, there are CNNs especially trained for place recognition that return either holistic image descriptors like NetVLAD \\cite{netvlad} or local features like DELF \\cite{delf}.\nThe performance of holistic descriptors can be further improved by additional pre- and postprocessing steps.\nTo improve the performance of holistic descriptors, feature standardization and unsupervised learning techniques like PCA- and clustering-based methods can be used \\cite{Schubert2020}.\nIn \\cite{Neubert2019} it is shown how a neurologically-inspired neural network can be used to combine several descriptors along a sequence to a new descriptor for each place.\n\\cite{Schubert2019} extended this approach to encode additional odometry information in the new descriptors.\nGiven the final pairwise similarities between descriptors from database and query, sequence-based methods \\cite{seqSLAM}\\cite{hmm} for postprocessing can be used to improve the place recognition performance further.\n\nIn this paper, we propose a graph-based approach to optimize the descriptor similarities by incorporating prior knowledge.\nIn \\cite{hmm} Hansen and Browning used a hidden Markov model (HMM) to formulate a graph-based sequence search method in the similarity matrix.\nNaseer et al. \\cite{Naseer14} used a flow network with edges defined for sequence search to improve place recognition results.\nIn \\cite{Zhang2019} a graph is used to prevent place matches between adjacent places and to distribute high matching scores to neighboring places.\nIn contrast to these approaches, our approach exploits not only sequence information but also integrates other additional knowledge, e.g., about intra-set similarities in the database set or the query set.\nThe literature provides several approaches for localization where the database is known in advance.\nFor example, given the images from the database (or another representative training set), FabMap \\cite{fabmap} learns statistics about feature occurrences in order to determine the most descriptive features.\nMcManus et al. \\cite{McManus2014} train offline condition-invariant broad-region detectors from beforehand collected images with a variety of appearances at particular locations.\nIn \\cite{Neubert2015} intra-database descriptor comparisons are used to reduce the number of required image comparisons during the query run.\nVysotska and Stachniss \\cite{Vysotska2017} use binary intra-database place matches to perform jumps within the similarity matrix during sequence search.\nIn contrast, in our presented approach we use potentially continuous intra-database similarities to optimize the place recognition result instead of just accelerating it.\nMoreover, we do not only use this information to find loops but in addition to potentially inhibit wrong loop closures.\n\nGraphical models, and in particular factor-graphs, are a well established tool in mobile robotics \\cite{Dellaert2017}, e.g., in the form of robust pose graph SLAM \\cite{Sunderhauf2012}.\nSimilar to the here proposed approach, pose graph SLAM represents each place as a node and the edges (factors) model constraints between these places. \nHowever, a significant difference is that pose graph SLAM deals with spatial information, i.e., the places are represented by pose coordinates in the world and the factors are spatial transformations between these poses (e.g., odometry or detected loop closures). \nIn contrast, our presented approach is intended to be used \\textit{before} SLAM to establish loop closures. \nIn particular, we do not optimize metric errors between spatial constraints but errors in the mutual consistency of descriptor similarities.\n\n\n\\section{ARCHITECTURE OF THE GRAPHICAL MODEL}\\label{sec:algo}\nA graphical model serves as a structured representation of prior knowledge in terms of dependencies, rules, and available information.\nGiven the variable nodes in the graph with their initial values together with dependencies between nodes based on additional knowledge, optimization algorithms can be used to modify the variables in order to resolve inconsistencies between nodes.\nWe are going to exploit this mechanism and present a graph-based framework for visual place recognition that optimizes the initial similarity values $\\hat{S}^{DB\\times Q}$ from pairwise image descriptor comparisons.\n\nThe graph-based framework is expressed as factor graph.\nFactor graphs are graphical representations of least squares problems -- for a detailed introduction to factor graphs please refer to \\cite{Dellaert2017}.\nThe graph's basic architecture consists of nodes with unary factors that penalize deviations from the initial similarity values.\nDepending on additional knowledge, we can add different factors in the graph to introduce connections (i.e. dependencies) between nodes.\nTwo architectures with nodes, unary factors, and additional factors are shown in Fig.~\\ref{fig:graph_intra_set}. \nEach factor defines a quadratic cost function based on the variables it connects. \nThe resulting combined optimization problem is defined as a weighted sum of the individual cost functions.\nThe optimization is subject of Sec.~\\ref{sec:optim}.\nHere, we first define the basic architecture of the graph with corresponding nodes and unary factors.\nNext, we propose two ways to exploit prior knowledge with binary and n-ary factors.\nWe will structure the explanation of each factor by the \\textit{prior knowledge} that can be exploited, a corresponding \\textit{rule} that expresses the resulting dependency between nodes, a proposed related \\textit{cost function} $f(\\cdot)$ that punishes a violation of this rule, and the factor's \\textit{usage} in the graph.\nExcept for the unary factor, each factor is optional and depends on the available knowledge.\nAccordingly, the proposed framework can be extended in future work with additional factors.\n\n\\subsection{Graph nodes}\\label{sec:algo:nodes}\nThe graph-based framework is designed to optimize the initial pairwise descriptor similaritiy matrix $\\hat{S}^{DB\\times Q} \\in \\mathbb{R}^{M\\times N}$.\n$M$ is the number of database images and $N$ the number of query images.\nAccordingly, we define $M\\times N$ nodes where each node $s_{ij} \\in S^{DB\\times Q}$ is a variable version of its initial value $\\hat{s}_{ij} \\in \\hat{S}^{DB\\times Q}$ with $s_{ij},\\hat{s}_{ij} \\in [0,1]$.\n\n\\subsection{Unary factor} \\label{sec:unary_factor}\n\\subsubsection{Prior Knowledge}\nDescriptor similarities $\\hat{s}_{ij}$ are the primary source of information for place recognition.\nBeyond initializing the variables to these similarities, we have to prevent too large deviations of $s_{ij}$ from $\\hat{s}_{ij}$, in particular to prevent trivial solutions during optimization.\n\n\\subsubsection{Rule ''prior``}\n\\begin{align}\\label{eq:rule_unary}\n s_{ij} \\approx \\hat{s}_{ij}\n\\end{align}\n\n\\subsubsection{Cost function ''prior``}\n\\begin{align}\\label{eq:cost_unary}\n f_\\text{prior}(\\cdot) = (s_{ij} - \\hat{s}_{ij})^2\n\\end{align}\n\n\\subsubsection{Usage}\nEach node $s_{ij}$ is connected with a single unary factor to its initial value $\\hat{s}_{ij}$ as shown in Fig.~\\ref{fig:graph_intra_set}.\nThus, $M\\times N$ unary factors are used in a graph.\n\n\\subsection{Binary factors for the exploitation of intra-database or intra-query similarities from poses or descriptor-comparisons}\\label{sec:binary_factor}\n\n\\subsubsection{Prior Knowledge -- intra-database similarities from poses}\nIn some applications, the database is created with a more advanced sensor setup than the query.\nDue to the missing position information for the query images, direct position-based matching cannot be conducted.\nNevertheless, available poses for database images can be used to create binary intra-database similarities $\\hat{s}^{DB}_{ij} \\in \\hat{S}^{DB}$ that encode whether two database images $i$ and $j$ show the same place ($\\hat{s}^{DB}_{ij}:=1$) or different places ($\\hat{s}^{DB}_{ij}:=0$).\n\n\\subsubsection{Prior Knowledge -- intra-set similarities from image descriptors}\nEven if position data is not available, we can acquire similar information directly from image comparisons within the database and also within the query to get intra-database similarities $\\hat{s}^{DB}_{ij} \\in \\hat{S}^{DB}$ and intra-query similarities $\\hat{s}^Q_{ij} \\in \\hat{S}^Q$ with $\\hat{s}^{DB}_{ij}, \\hat{s}^Q_{ij} \\in [0,1]$.\nThese intra-set similarities are an inherent and almost always available source of additional information, which has not been used often yet.\nThe computation of intra-set similarities can usually be done more reliably than the computation of inter-set similarities, because the condition \\textit{within} database or query is potentially more constant than between both. \nThe intra-set image comparisons could be done with methods like FabMap~\\cite{fabmap} that are more suited for place recognition under constant condition.\nFig.~\\ref{fig:s_matrices} shows how the intra-set similarities $\\hat{S}^{DB}$, $\\hat{S}^Q$ are related to the inter-set similarities $\\hat{S}^{DB\\times Q}$ and how they can reveal loops and zero-velocity stages.\n\n\\begin{figure}[tb]\n \\centering\n \\includegraphics[width=0.8\\linewidth]{images\/S_matrices.pdf}\n \\vspace{-0.2cm}\n \\caption{\n The information of the similarity matrix between database and query (bottom right) can be extended with similarity information of images \\textit{within} the query set (top right) and\/or \\textit{within} the database set (bottom left).}\n \\vspace{-0.2cm}\n \\label{fig:s_matrices}\n\\end{figure}\n\nWe define separate binary factors for both intra-set similarities, either from poses or from image descriptors.\nEach binary factor is a combination of two complementing rules and cost functions $f^Q_\\text{loop}(\\cdot)+f^Q_\\text{exclusion}(\\cdot)$ (Eq.~(\\ref{eq:cost_sim_q}), (\\ref{eq:cost_dissim_q})) and $f^{DB}_\\text{loop}(\\cdot)+f^{DB}_\\text{exclusion}(\\cdot)$ (Eq.~(\\ref{eq:cost_sim_db}), (\\ref{eq:cost_dissim_db})), respectively.\n\n\\subsubsection{Rule ''loop``}\nIf $\\hat{s}^{Q}_{jl}$ is high (indicated by ''$\\hat{s}^{Q}_{jl}\\uparrow$``), the $j$-th and $l$-th query image likely show the same place.\nIf the $i$-th database image is compared to these two query images, the following ternary relation can be derived:\n\\begin{align}\\label{eq:rule_sim2}\n s_{ij} \\approx s_{il} &\\text{, iff } \\hat{s}^Q_{jl}\\uparrow\n\\end{align}\n\nThe rule expresses that ''$s_{ij}$ should be equal to $s_{il}$ iff the query images $j$, $l$ show the same place``, because if both query images $j$, $l$ show the same place, database image $i$ can either be equal to both or to none of both.\nAccordingly, this rule exploits loops within the intra-query similarities.\nThis rule is inherent to the place recognition problem and valid as well for intra-database similarities:\n\\begin{align}\\label{eq:rule_sim1}\n s_{ij} \\approx s_{kj} &\\text{, iff } \\hat{s}^{DB}_{ik}\\uparrow\n\\end{align}\n\n\\subsubsection{Cost function ''loop``}\nFor equation (\\ref{eq:rule_sim2}) and (\\ref{eq:rule_sim1}), cost functions similar to (\\ref{eq:cost_unary}) can be formulated:\n\\begin{align}\n f^Q_\\text{loop}(\\cdot) &= \\hat{s}^Q_{jl} \\cdot (s_{ij} - s_{il})^2 \\label{eq:cost_sim_q} \\\\\n f^{DB}_\\text{loop}(\\cdot) &= \\hat{s}^{DB}_{ik} \\cdot (s_{ij} - s_{kj})^2 \\label{eq:cost_sim_db} \n\\end{align}\nThe quadratic error term is multiplied by the intra-set similarity $\\hat{s}^Q_{jl}$ or $\\hat{s}^{DB}_{ik}$ to apply weighting in case of non-binary intra-set similarities from image descriptors or to turn it on and off for binary intra-set similarities.\n\n\\subsubsection{Rule ''exclusion``}\nIf $\\hat{s}^{Q}_{jl}$ is low (indicated by ''$\\hat{s}^{Q}_{jl}\\downarrow$``), the $j$-th and $l$-th query image probably show different places.\nIf the $i$-th database image is compared to these two query images, the following ternary relation can be derived:\n\\begin{align}\\label{eq:rule_dissim2}\n \\neg(s_{ij}\\uparrow \\land s_{il}\\uparrow) \n &\\text{, iff } \\hat{s}^Q_{jl}\\downarrow\n\\end{align}\nThis rule expresses that ''not both similarity measurements $s_{ij}$ AND $s_{il}$ can be high iff the query images $j$, $l$ show different places``, i.e., the rule excludes one or both similarities $s_{ij}$, $s_{il}$ from being high; otherwise a single database image $i$ would show two different places concurrently.\nAgain, this rule is inherent to the place recognition problem and is supposed to add valuable information that can be exploited.\nAs before, the same rule is valid for intra-database similarities:\n\\begin{align}\\label{eq:rule_dissim1}\n \\neg(s_{ij}\\uparrow \\land s_{kj}\\uparrow) \n &\\text{, iff } \\hat{s}^{DB}_{ik}\\downarrow\n\\end{align}\n\n\\subsubsection{Cost function ''exclusion``}\nIt is less natural how to express the \\textit{rule ''exclusion``} in a cost function. This is part of the discussion in Sec.~\\ref{sec:discussion}.\nIn this work, we define the following cost functions for equation (\\ref{eq:rule_dissim2}) and (\\ref{eq:rule_dissim1}):\n\\begin{align}\n f^Q_\\text{exclusion}(\\cdot) &= (1-\\hat{s}^Q_{jl}) \\cdot (s_{ij} \\cdot s_{il})^2 \\label{eq:cost_dissim_q} \\\\\n f^{DB}_\\text{exclusion}(\\cdot) &= (1-\\hat{s}^{DB}_{ik}) \\cdot (s_{ij} \\cdot s_{kj})^2 \\label{eq:cost_dissim_db}\n\\end{align}\nThe quadratic error term is multiplied by the negated intra-set similarity $(1-\\hat{s}^Q_{jl})$ or $(1-\\hat{s}^{DB}_{ik})$ to weight it in case of non-binary intra-set similarities from image descriptors or to turn it on and off for binary intra-set correspondences.\n\n\\subsubsection{Usage}\nFig.~\\ref{fig:graph_intra_set} (left) shows a graphical model with nodes $s_{ij}$, unary factors $f_\\text{prior}(\\cdot)$, and the proposed binary factors.\nTo apply the binary factors for existing intra-database similarities $\\hat{S}^{DB}$, each node $s_{ij}$ has to be connected to every node $s_{kj}$ for all $k = 1,\\ldots,M$ with $k\\neq i$ within each column in $S^{DB\\times Q}$; i.e., $\\binom{M}{2}$ factors per column.\nFor existing intra-query similarities $\\hat{S}^Q$, each node $s_{ij}$ has to be connected to every node $s_{il}$ for all $l = 1,\\ldots,N$ with $l\\neq j$ within each row in $S^{DB\\times Q}$; i.e., $\\binom{N}{2}$ factors per row.\nThe potentially high number of connections is part of the discussion in Sec.~\\ref{sec:discussion}.\n\n\\begin{figure}[tb]\n \\centering\n \\includegraphics[width=\\linewidth]{images\/graph.pdf}\n \\caption{Illustration of the graph structure. \\textit{(left)} Unary and binary factors. \\textit{(right)} Unary and n-nary factors, which connect structured local blocks.}\n \\label{fig:graph_intra_set}\n\\end{figure}\n\n\\subsection{N-ary factors for the exploitation of sequences}\\label{sec:nary_factor}\n\n\\subsubsection{Prior Knowledge}\nIf both database and query are recorded as spatio-temporal sequence, sequences also appear in the inter-set similarities $S^{DB\\times Q}$ (see Fig.~\\ref{fig:s_matrices}).\nIn this case, SeqSLAM \\cite{seqSLAM} showed the benefit from a simple combination of similarities of neighboring images.\nOriginally, it was implemented as postprocessing of the similarity matrix $S^{DB\\times Q}$. \nWe can integrate such sequential information within the graph in a similar fashion using n-ary factors.\n\n\\subsubsection{Rule ''sequence``}\n\\begin{align}\\label{eq:rule_seq}\n s_{ij} \\approx \\frac{1}{L}\\sum_{\\forall \\{k,l\\} \\in \\text{Seq}({i,j};L,v_p)} s_{kl}\n\\end{align}\n$\\text{Seq}({i,j};L,v_p)$ is a function that returns similarities that are part of a sequence along a line segment with slope $v_p$ and sequence length $L$ within $S^{DB\\times Q}$ around $s_{ij}$.\nFor a full explanation how SeqSLAM works, please refer to \\cite{seqSLAM}.\n\n\\subsubsection{Cost function ''sequence``}\n\n\\begin{align}\n f_\\text{seq}(\\cdot) = (s_{ij} - \\frac{1}{L}\\max_{v_p \\in V}{(\\sum_{\\{k,l\\} \\in \\text{Seq}(\\{i,j\\};L,v_p)} s_{kl})})^2\n\\end{align}\nwith $V$ being the set of allowed velocities within $S^{DB\\times Q}$.\n\n\\subsubsection{Usage}\nFig.~\\ref{fig:graph_intra_set} (right) shows a graphical model with nodes $s_{ij}$, unary factors $f_\\text{prior}(\\cdot)$, and the proposed n-ary factors for sequences.\nAs for the unary factors in Sec.~\\ref{sec:unary_factor}, each node in the graph is equipped with one n-ary factor.\nEach n-ary factor connects its node to all nodes that are part of any sequence with length $L$ and slope $v_p\\in V$.\nAccordingly, $M\\times N$ n-ary factors are introduced into a graph if sequences are exploited.\n\n\\section{OPTIMIZATION OF THE GRAPHICAL MODEL}\\label{sec:optim}\nThe objective of the graph-based framework is the formulation of dependencies from prior knowledge for a subsequent optimization of the similarities $s_{ij}$ to incorporate the prior knowledge into the final similarity values, and to potentially resolve contradictory dependencies.\n\nWe use factor graphs as graphical representation of least squares problems and defined every cost function for each type of factor as (weighted) quadratic error function.\nAccordingly, the global error $E$ is defined as a sum over the (weighted) cost functions of all factors in the graph.\n\n\\subsection{Weighting of factor's costs for global error computation in the graph}\nAs usually done in error computation for a graph, costs from different factor-types have to be balanced by weighting them separately.\nTherefore, we normalize the cost of each factor by the number of factors per factor-type, and introduce user-specified weights $w$ for all factors except for the basic unary factor (which gets weight 1).\n\nFor the unary factor of our basic architecture with cost function $f_\\text{prior}(\\cdot)$, we define the partial global error $E_\\text{prior}$ with\n\\begin{align}\n E_\\text{prior} = \\frac{1}{MN} \\sum_{i=1}^M \\sum_{j=1}^N f_\\text{prior}(s_{ij})\n\\end{align}\n\nIn case of dependencies in the graph that are introduced by available intra-database or intra-query similarities, we define a partial global error $E^{DB}_{loop,exclusion}$ or $E^Q_{loop,exclusion}$ for the weighted summation over all binary factors (Sec.~\\ref{sec:binary_factor}): \n\\begin{align}\n E^{DB}_\\text{loop,exclusion} = &\\frac{1}{N\\binom{M}{2}} \\sum_{j=1}^N \\sum_{i=1}^{M-1} \\sum_{k=i+1}^M w^{DB}_\\text{loop} \\cdot f^{DB}_\\text{loop}(s_{ij}, s_{kj}) \\nonumber\\\\\n &+ w^{DB}_\\text{exclusion} \\cdot f^{DB}_\\text{exclusion}(s_{ij}, s_{kj})\n \\\\\n E^Q_\\text{loop,exclusion} = &\\frac{1}{M\\binom{N}{2}} \\sum_{i=1}^M \\sum_{j=1}^{N-1} \\sum_{l=j+1}^N w^Q_\\text{loop} \\cdot f^Q_\\text{loop}(s_{ij}, s_{il}) \\nonumber\\\\\n &+ w^Q_\\text{exclusion} \\cdot f^Q_\\text{exclusion}(s_{ij}, s_{il})\n\\end{align}\n\nFor the n-ary factors (Sec.~\\ref{sec:nary_factor}) in case of sequence exploitation, the partial global error $E_\\text{seq}$ is defined with\n\\begin{align}\n E_\\text{seq} = \\frac{1}{MN} \\sum_{i=1}^M \\sum_{j=1}^N w_\\text{seq} \\cdot f_\\text{seq}(s_{ij}, \\text{Seq}(\\{i,j\\};L,v^*_{ij})\n\\end{align}\nwith $v^*_{ij}$ being the optimal velocity (slope) with the highest mean of connected similarities around $s_{ij}$.\n\nFinally, summation over all partial global errors $E_i$ that occur in the graph yields the global error $E$:\n\\begin{align}\n E = \\sum_{\\forall i} E_i\n\\end{align}\n\n\n\\subsection{Implementation of the optimization}\nError $E$ is a sum solely over quadratic cost functions.\nThus, many tools for non-linear least squares (NLSQ) optimization can be used (e.g., scipy's \\textit{least\\_squares}-function in Python).\nFor an easier formulation of the optimization problem, factor graph tools like g2o \\cite{g2o} can be used.\nHowever, one should be aware that depending on the introduced factors, the optimization problem can get quite huge.\nThus, efficient implementations should be preferred that perform a quick and memory efficient optimization.\nSec.~\\ref{sec:seq} reports our achieved runtimes that are presumably sufficient for many applications.\nSec.~\\ref{sec:discussion} provides some more discussion on alternative optimization methods and approximation techniques for runtime improvements.\n\n\\section{EXPERIMENTAL RESULTS}\\label{sec:results}\nWe present experiments that investigate the performance gains achieved by the graph-based framework with three different extensions that exploit 1)~intra-database similarities; 2)~intra-database and intra-query similarities; 3)~intra-database similarities, intra-query similarities and sequences.\nIn order to evaluate the potential benefit beyond available pre- and postprocessing methods, we repeat these experiments with the descriptor standardization approach from \\cite{Schubert2020} for preprocessing and SeqSLAM for postprocessing.\nIn a final experiment, we compare our graph-based method with sequence-based approaches from the literature.\n\n\\subsection{Experimental Setup}\n\\subsubsection{Image descriptor}\nNetVLAD \\cite{netvlad} is used as CNN-image descriptor in all experiments.\nWe use the author's implementation trained on the Pitts30k dataset with VGG-16 and whitening.\n\n\\subsubsection{Metric}\nThe performance is measured with average precision which is the area under the precision-recall curve.\n\n\\subsubsection{Datasets}\nAll experiments are based on the five different datasets Nordland~\\cite{nordland}, StLucia (Various Times of the Day)~\\cite{stlucia}, CMU (Visual Localization)~\\cite{cmu}, GardensPoint Walking~\\cite{gardenspoint} and Oxford RobotCar~\\cite{robotcar} with different characteristics regarding the environment, appearance changes, in-sequence loops, stops, or viewpoint changes.\nWe use the datasets as described in our previous work~\\cite{Neubert2019}.\nImages from StLucia, CMU and Oxford were sampled with one frame per second, which preserves varying camera speeds, stops and loops in the datasets, and leads to translation and orientation changes during revisits. GardensPoint contains lateral viewpoint changes.\n\nFor the binary intra-database similarities $\\hat{S}^{DB}_{pose}$ from poses, we use the GPS from the datasets or a main diagonal for \\textit{GardensPoint Walking} and \\textit{Nordland}.\nFor the query sequence, we assume that no GPS is available.\n\n\\subsubsection{Implementation}\\label{sec:implementation}\nGraph creation and optimization were implemented in Python with scipy's \\textit{least\\_squares}-optimization function; the \\textit{Trust Region Reflective} algorithm was used for minimization.\nDue to the huge number of factors within a graph in case of the usage of intra-database and intra-query similarities, we divided $S^{DB\\times Q}$ into equally sized patches with height and width $\\leq 500$, and optimized each patch separately.\nNo information is shared between patches, and the n-ary factors are truncated at borders.\nA full optimization on $S^{DB\\times Q}$ without patches was performed if only intra-database similarities $\\hat{S}^{DB}$ were used.\nThe variables in the optimization are initialized with $\\hat{S}^{DB\\times Q}$ from the pairwise descriptor comparisons.\n\n\\subsubsection{Parameters}\nIn all experiments, we used a fixed parameter set that was determined from a toy-example and a small real-world dataset.\nWe used $w^{DB}_\\text{loop}=4, w^{DB}_\\text{exclusion}=40$ for intra-set similarities from GPS, $w^{DB}_\\text{loop}=w^Q_\\text{loop}=1, w^{DB}_\\text{exclusion}=w^Q_\\text{exclusion}=20$ for intra-set similarities from descriptors, and $w_\\text{seq}=10$, $L=11$ for sequences.\n\n\\subsection{Contributions of information sources and rules}\\label{sec:exp1}\nIn Sec.~\\ref{sec:algo} we identified four rules: \\textit{''prior``}, \\textit{''loop``}, \\textit{''exclusion``} and \\textit{''sequence``}.\nIn the following, we are going to evaluate the influence of the rules when they are successively added and exploited in the graph.\nNote that rule \\textit{''prior``} alone would merely return the initial values $\\hat{S}^{DB\\times Q}$.\nInput to the graph are the pairwise descriptor similarities from the raw NetVLAD descriptors that serve as baseline as well (termed ''pairwise``).\nAll results are summarized in Table~\\ref{tab:no_std_seq}.\n\\input{tables\/res.tex}\n\n\n\\subsubsection{Exploitation of intra-database similarities from poses or descriptors (rules \\textit{''loop``} and \\textit{''exclusion``})}\\label{sec:exp2}\nIntra-database similarities $\\hat{S}^{DB}$ can be used in most place recognition setups as they can be acquired either from the pairwise image comparisons within the database or from poses, e.g., from GPS or SLAM.\n\nTable~\\ref{tab:no_std_seq} shows the results (indicated by $\\hat{S}^{DB}_{pose}$ and $\\hat{S}^{DB}_{desc}$) when intra-database similarities are used either from poses or descriptor comparisons.\nIn most cases, the pairwise performance is significantly improved and never gets worse.\nIf intra-database similarities from poses (here: GPS) are used, the performance gain is slightly better, presumably because place matchings and distinctions from poses are binary and less error prone.\nHowever, even when the intra-database similarities from descriptor comparisons are used, the performance can be improved significantly.\n\nThe results support that most existing place recognition pipelines could be improved with the proposed graph-based framework together with intra-database similarities.\nMoreover, since it is not necessary to know the query sequence in advance, the proposed graph-based framework can be employed in an online fashion.\n\n\n\\subsubsection{Exploitation of intra-database and intra-query similarities from poses or descriptors (rules \\textit{''loop``} and \\textit{''exclusion``})}\nIn addition to the previous experiment, supplementary intra-query similarities could be used to model dependencies within the graph not only between database images but also between query images.\nWe used intra-query similarities solely from descriptor comparisons since we do not assume global pose information during the query run; otherwise, place matchings could be received directly from pose-comparisons.\n\nTable~\\ref{tab:no_std_seq} shows the results (indicated by $\\hat{S}^{DB}_{pose}{+}\\hat{S}^Q_{desc}$ and $\\hat{S}^{DB}_{desc}{+}\\hat{S}^Q_{desc}$).\nFor intra-database similarities from poses, the performance is improved only for a few sequence-combinations compared to the performance with intra-database but without intra-query similarities.\nWhen using intra-query similarities in addition to intra-database similarities from descriptors, the performance could be improved at least by additional $5\\%$ for $50\\%$ of all datasets.\n\nThe results indicate that additional data from intra-query similarities within the proposed graph-based framework can be used to improve the place recognition performance further.\nIt is important to note that the system again never performs worse in comparison to a graph with intra-database but without intra-query similarity exploitation.\nThe result is interesting, since intra-query similarities from descriptors could always be collected for a subsequent postprocessing.\n\n\\subsubsection{Additional exploitation of sequences within the graph (rules \\textit{''loop``}, \\textit{''exclusion``} and \\textit{''sequence``})}\nIn this experiment, we used sequence information within the graph in addition to the intra-database and intra-query similarities.\nTable~\\ref{tab:no_std_seq} shows the results (indicated by $\\hat{S}^{DB}_{pose}{+}\\hat{S}^Q_{desc}{+}\\text{Seq}$ and $\\hat{S}^{DB}_{desc}{+}\\hat{S}^Q_{desc}{+}\\text{Seq}$).\nWith sequence information the graph can again significantly improve the place recognition performance compared to the previous experiments without this additional assumption no matter if the intra-database similarities come from poses or descriptor comparisons.\nMoreover, the full setup of the graph with all proposed factors clearly outperforms the baseline.\n\n\\subsection{Combination with preprocessed descriptors}\\label{sec:exp2_std}\nThe place recognition performance can be improved if the descriptors are preprocessed \\cite{Schubert2020}.\nTo investigate the influence of preprocessing on all configurations of the graph from Sec.~\\ref{sec:exp1}, we repeated all previous experiments with standardized image descriptors \\cite{Schubert2020}.\nResults are shown in Table~\\ref{tab:std_seq} (left).\nSince the preprocessing requires complete knowledge of the query descriptors, we do not provide results for the online configuration of the graph.\nThe full configuration of the graph with intra-set similarities and sequences can again show significant performance improvements for the majority of the sequence-combinations.\nThe results demonstrate that the graph-based framework benefits from better performing descriptors.\n\\input{tables\/res_cat_STD_seq.tex}\n\n\\subsection{Combination with sequence-based postprocessing}\\label{sec:exp_postproc}\nIn the next experiment, a modified version of SeqSLAM~\\cite{seqSLAM} without local contrast normalization and the single matching constraint is used; otherwise SeqSLAM would fail on datasets with in-sequence loops.\nIt postprocesses the inter-set similarities $S^{DB\\times Q}$ either from the pairwise descriptor comparison (baseline) or from all configurations of the graph from Sec.~\\ref{sec:exp1}.\n\nResults are shown in Table~\\ref{tab:std_seq} (right);\nnote that descriptor comparisons from the raw NetVLAD descriptors \\textit{without} sequence-based postprocessing are used as inputs to the graphs.\nThe baseline performance after postprocessing compared to the baseline performance without postprocessing (Table~\\ref{tab:no_std_seq}) is already comparatively high.\nAccordingly, it is hard to achieve high performance improvements.\nNonetheless, the graph-based framework could improve the results in almost $50\\%$ of the cases often by more than $10\\%$ -- however, intra-database similarities from both poses and descriptors seem to be sufficient, since more information could not be used for further performance improvements on most datasets.\n\n\\subsection{Combination with pre- and postprocessing}\nWe also conducted experiments with preprocessing from Sec.~\\ref{sec:exp2_std} \\textit{and} postprocessing from Sec.~\\ref{sec:exp_postproc}.\nThe baseline performance got already almost perfect for most of the datasets, so that performances could only be improved slightly by less than $5\\%$ with the graph-based approach.\nAgain, the performance was never worse than the baseline.\n\n\\subsection{Comparison with state-of-the-art sequence-based methods}\\label{sec:seq}\nTo compare performance and runtime of our method with approaches from the literature, we conduced additional experiments with the sequence-based methods SeqSLAM~\\cite{seqSLAM}, MCN~\\cite{Neubert2019}, VPR~\\cite{Vysotska2017} and ABLE~\\cite{able}.\nThe sequence length was $L=11$, if required.\nThe experiments were performed twice without and with feature standardization~\\cite{Schubert2020} for descriptor preprocessing.\n\nTable~\\ref{tab:seq} shows the achieved performances.\nOur method clearly outperformed the compared approaches and achieved the best performance on most datasets.\nOnly on Nordland SeqSLAM and ABLE achieved better performance, since they benefit from constant camera speed in database and query.\nWith feature standardization, ABLE could additionally achieve best performance on GardensPoint, and MCN performed best on three Oxford datasets.\n\nRuntimes were measured on an Intel i7-7700K CPU with 64GB RAM.\nThe maximum runtimes per query for all methods are shown in Table~\\ref{tab:seq} (bottom).\nOur method required approx. 5.7sec per query on Oxford\\#2 with 3413 database images, while all other approaches needed less than 500msec.\nAll reported runtimes are presumably sufficient for applications like loop closure detection in SLAM, since loops need not to be detected with the full frame rate of a camera.\nSec.~\\ref{sec:discussion} provides some more discussion on the computational efficiency and potential improvements.\n\n\\input{tables\/res_seq.tex}\n\n\\section{DISCUSSION AND CONCLUSION}\\label{sec:discussion}\nThe previous sections presented our approach to use a graphical model as a flexible framework to model different kinds of additional information available in place recognition.\nThe experiments demonstrated that the presented method can considerably improve place recognition results in various configurations in terms of available data (e.g., poses of database images), the subset of applied rules (e.g., using sequences or not), or restriction to online place recognition.\nThe representation and optimization using graphical models offers a high degree of flexibility.\nIn the remainder of this last section, we want to discuss aspects of the proposed system and some particularly interesting possible extensions.\n\nThe graph-based optimization is performed on the pairwise descriptor similarities $S^{DB\\times Q}$.\nThis makes it relatively independent of the actually chosen image descriptor, which may require only a slight parameter adjustment.\nThis property even allows an optimization of place descriptor similarities from different sensor modalities like LiDAR.\n\nIn this paper, we defined several factor-types to express prior knowledge (``rules'') about place recognition problems.\nEach factor implements a cost function for a rule. \nOften, there are alternative formulations of a cost function for a rule.\nFor example, the cost function in case of no loop in query ($\\hat{s}^Q_{ij}\\downarrow$; Eq.~(\\ref{eq:cost_dissim_q})) or database ($\\hat{s}^{DB}_{ij}\\downarrow$; Eq.~(\\ref{eq:cost_dissim_db})) is defined as multiplication $(s_1\\cdot s_2)^2$.\nEspecially for this rule, alternative formulations may apply like $\\min(s_1, s_2)^2$ or $\\max(s_1+s_2-1, 0)^2$; both are piecewise linear which could be beneficial.\n\nFactor graphs are often (but not exclusively) used in combination with probabilities. Presumably, a more probabilistic view on the proposed graph-based framework could provide additional insights.\nFor instance, the chosen cost functions (\\ref{eq:cost_unary}), (\\ref{eq:cost_sim_db}) and (\\ref{eq:cost_dissim_q}) with structure $(s_1-s_2)^2$ can be considered as the negative log-likelihood of a single Gaussian:\n\\vspace{-0.02cm}\n\\begin{align}\n (s_1-s_2)^2 \\Leftrightarrow -\\ln(e^{-(s_1-s_2)^2})\n\\end{align}\nMoreover, a piecewise linear formulation of the discussed factor above could help to formulate the proposed graph-based framework in a more probabilistic way; for instance a cost function $\\min(s_1, s_2)^2$ corresponds to the negative log-likelihood of a maximum of two Gaussians:\n\\begin{align}\n \\min(s_1, s_2)^2 \\Leftrightarrow -\\ln(\\max(e^{-(s_1-0)^2}, e^{-(s_2-0)^2})\n\\end{align}\n\nThese cost functions, however, solely work on the descriptor similarities.\nIn the related work (Sec.~\\ref{sec:related_work}), we already mentioned the important difference to pose graph SLAM, which solely works on spatial poses (and not their similarities).\nAn interesting question for future work is whether the proposed approach can be extended to also directly work on descriptors instead of their similarities (i.e., the variables would be descriptor vectors, not their scalar similarities).\nThis significantly increases the complexity of the optimization problem, but could allow the simultaneous optimization of spatial poses and descriptors for a potentially tightly coupled loop closure detection and SLAM.\n\nEven without such an extension, as indicated in Sec.~\\ref{sec:algo} and Sec.~\\ref{sec:optim}, the number of factors or connections between nodes can get quite huge, and grows cubically if intra-set similarities are used.\nIn our experiments, we addressed the problem by dividing $S^{DB\\times Q}$ into patches if both intra-set similarities were used (Sec.~\\ref{sec:implementation}).\nPresumably, using more efficient implementations, e.g. C++ implementations in Ceres \\cite{ceres}, can improve memory consumption and computational efficiency.\nAnother promising direction are approximation techniques like a systematic removal of low-relevance connections in the graph.\nFinally, different optimization techniques could be used: In earlier work on graph optimization, minimization techniques based on hill-climbing algorithms like ICM (iterated conditional modes) were used \\cite[p.599]{Koller2009}; these may allow a different and more compact representation and optimization of the graph.\n\n\\input{root.bbl}\n\n\\end{document}\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\nThe ubiquity of opportunity offered by the Internet of Things (IoT) is providing new ways to embed computation and storage into everyday devices. Reductions in the cost of hardware, coupled with technological advancements, are enabling a constant stream of innovative new products and services that rely upon IoT architectures. Such developments are particularly evident in the `wearable' and healthcare industries, with parallel innovations occurring in the Industrial Internet of Things (IIoT) \\cite{bessis2013,hill2017a}.\n\nAlong with the ability to sense environments, conditions and processes, comes the requirement to measure, monitor, analyse and evaluate performance. As such, Business Intelligence (BI) is undergoing a resurgence as consumer demand moves beyond traditional dashboards through forecasting, towards prescriptive analytics. Of late, vendors have delivered BI often as a flexible, on-demand service, making use of the underlying elastic cloud platforms that enterprise software applications are typically deployed upon \\cite{alaqrabi2012,alaqrabi2013,alaqrabi2015,alaqrabi2018}.\n\nAdvancements in technology have increased the ability to connect such a variety of embedded devices to larger pools of resources such as clouds. Integrating embedded devices and cloud servers raises an important discussion regarding the nature of data generated or transmitted by IoT devices. These approaches must be secure and provide the necessary privacy controls for users. At present, the security and privacy concerns created by these devices play a central role in the successful integration of these two technologies \\cite{kalra2015}.\n\nThe heterogeneous nature of IoT environments makes it much harder to detect the insider and outsider attacks in such universal platforms \\cite{Alrawais2017}.\n\nExperiences with clouds, especially those that have public or hybrid architectures, illustrates that cloud services and applications require multi-layered approaches to external and internal threats to security. This is most pertinent in the IIoT environment, where business systems contain valuable Intellectual Property (IP) that can only be protected by retaining tight security over working practices.\n\nThis article describes a model for delivering analytics services to and between components such as those encountered in IoT architectures. A number of adversarial attacks are simulated to demonstrate the effectiveness of a cloud-inspired multi-layer security model in the IoT domain.\n\\section{Cloud Computing Model}\nThe core principles of cloud computing: on-demand self-service, broad network access, resource pooling, rapid elasticity and service metering (NIST Special Publication 800-145) \\cite{hill2013}, make utility computing an attractive proposition for environments that contain distributed resources.\nThe NIST definition \\cite{nist2011} describes five options for the deployment of cloud computing as follows: public clouds, private clouds hosted onsite, private clouds hosted off-site, onsite community clouds and offsite community clouds.\n\nAbstraction is one of the most compelling motivations for system architects; this enables myriad heterogeneous resources to be viewed as a homogeneous whole. Cloud architectures of any of the five types defined by NIST are modelled using a framework that consists of seven layers. The physical infrastructure foundation is layer one, upon which a resource abstraction (virtualisation) second layer resides. Layer three is the resource composition layer, and layer four refers to the Infrastructure as a Service model (IaaS). Platform as a Service (PaaS) exists on layer five, Software as a Service (SaaS) on layer seven and finally the cloud tenants' applications sit within layer seven \\cite{demchenko,hill2013}.\n\nWe have considered the application of analytics services in the context of an enterprise Business Intelligence application that is likely to be delivered across a heterogeneous network of devices. In a cloud environment this is typified by off-site private or community clouds, whereby multiple disparate business organisations can host their own enterprise software services and data repositories in what is, to all intents and purposes, a set of hardware resources that is dedicated to each tenant \\cite{demchenko2012,hill2013}.\nThe tenants access their dedicated resources via a VPN, thus maintaining clear separation between corporate services and their respective analytics implementations. However, should any of the businesses wish to share services to promote operational efficiencies, this can be enabled via community-based agreements \\cite{demchenko2012}.\n\nWe have used the cloud model as inspiration to propose a multi-layer service-oriented framework that can deliver secure services (such as BI analytics) across distributed infrastructure. Security and privacy are key concerns for tenants of shared cloud infrastructure such as that found in off-site private and community clouds. The sharing of services such as malware detection and inoculation has significant benefits for organisations who may not otherwise have strategies or the resources to maintain their own defences. We shall now briefly review services dedicated to security and privacy within clouds in the following section.\n\n\n\\section{Secure Service Delivery}\nIn keeping with the service orientation models described so far, security as a service can be considered as a multi-layer framework in itself \\cite{panian2008}, protecting each layer of the cloud computing model \\cite{carvakho2011}. Such a service lends itself to a utility offering in the same way that clouds are rapidly provisioned and expanded on demand, for a cost that is quantifiable and chargeable to the consumer.\n\nAs described above, resource abstraction through virtualisation is a key principle for a cloud-inspired architecture, and therefore any security and privacy services need to be made available to all relevant Virtual Machines (VM) for each client \\cite{kumar2011,luo2011}. \n\nTherefore, security and privacy control services must be orchestrated via appropriate service interfaces between each cloud and its respective tenants that are accessing the services \\cite{kumar2011}. This control will be managed by the relevant virtualisation security manager for each VM, and as a consequence the control must be located within the virtualisation layer.\n\nPrivacy is managed via policy resources that need to be retained securely in a protected space such as a digital vault for encryption keys or digital certificates and such like \\cite{pearson2009,diaz2013}, all contained within a tier above other layers such as authentication and client metadata layers \\cite{pervez2012,chadwick2012}.\n\nAny instances whereby session packets do not have the requisite authority to invoke a particular service interaction will result in the session being terminated, and is a key security feature that is embedded within the framework.\n\n\\section{Multi-layered Security Model}\nOur proposed model is illustrated in Figure \\ref{fig:multi} and describes the cloud-inspired architecture containing multiple layers. Each layer incorporates security and privacy services that make use of firewalls as gateways for session traffic from each of the respective clients.\n\nWe have implemented the model within Opnet, the description of which is as follows.\nEach firewall is represented by a Cisco PIX 535, which operates across the various layers including network, transport and application. Access control lists within the firewalls enable the governance of traffic based on IP addresses, protocols and ports for common and bespoke applications.\n\nThe firewalls are configured at the application layer to filter traffic from different URLs, encrypted sessions (HTTPS) and various clients such as Java, and are able to utilise Internetwork Key Exchange (IKE) to encrypt all sessions via DES, 3DES or AES.\n\nAs an example we have defined four separate LANs to represent different clients\/cloud tenants, and they each access the cloud through separate firewalls.\n\nTo illustrate an adversarial scenario, we have incorporated a simulated distributed attack from three malicious parties, who each are attempting to access the network through independent firewalls. Such an attack is a key concern for early adopters of IoT technologies, as the pervasive use of wireless communications presents many opportunities for business vulnerabilities to be exposed.\n\n\\begin{figure}[tb]\n \\includegraphics[width=\\linewidth]{.\/figures\/Multilayer_hierarchical_Inter_Cloud_architecture.JPG}\n \\caption{Multilayer hierarchical inter-cloud architecture.}\n \\label{fig:multi}\n\\end{figure}\nAll of the multiple layers of the cloud are illustrated in Figure \\ref{fig:cloudlay}. We have embraced the cloud concept of resource abstraction and exploited this in the multi-layer model by representing each of the layers as separate clouds. Within each cloud layer exists an array of computing hardware to maximise performance and therefore minimise response times to service requests. It is an imperative that the model must not introduce excessive overheads into the normal functionality of the system.\n\nThe cloud layers are arranged such that inbound traffic from clients are filtered by the firewalls and then passed through successive layers until the system is satisfied that the requests can be delivered to the analytics functionality residing on the cloud applications layer. We now describe each of the layers and the role that they play within the model in turn:\n\n\\begin{itemize}\n\\item {\\bf Tenant firewalls}. This cloud layer consists of a number of databases that hold authentication data for each of the tenants, in order for them to be permitted access to the VM that have been assigned to them as part of their subscription. This is the gateway for each session, $S$ to be invoked.\n\n\\item {\\bf Tenant metadata}. Beyond the authentication criteria that is marshalled by the tenant firewalls layer, there exists further metadata for each client tenant. This metadata describes the credentials to authorise access to specific instances of repositories, applications and services. The detail is embedded within the session packets and is used to verify whether the session can continue or not.\n\\item {\\bf Digital vaults}. The vaults are a secure place to retain digital signatures\/certificates and decryption keys so that only the requisite authorities can access their own content and services. Public key encryption ensures that legitimate session requests are honoured and bogus requests are terminated.\n\\item {\\bf Intrusion Prevention System}. The occurrence of adversarial attacks is not limited to the correct authentication at the start of a session. Intrusion detection prevents attacks on sessions that are in progress, such as SQL injection in web forms. This might be considered a route into the DB\\textsubscript{META} or DB\\textsubscript{VAULT} repositories from the perspective of an adversary. This cloud layer prevents such activity from continuing.\n\n\\item {\\bf Malware protection}. Similar to the IPS cloud layer, an anti-malware layer contiuously montitors for trojan activity, where malicious exploits are embedded and concealed within session packets, only to be executed at the application layer. Records of the monitoring and detection history are retained within DB\\textsubscript{ANTIMAL}.\n\\item {\\bf Tenant applications}. This layer hosts the tenant's applications themselves, in this case the analytics functionality of enterprise BI. The preceding layers ensure that only a marshalled session, $S$ can access this layer, which in the case of BI potentially provides access to confidential business operations and performance data. For some functionality, there will be a requirement for further authentication from a user, such as an account and password, etc.\n\n\\end{itemize}\n\\subsection{Tenant repositories}\nThe final cloud layer hosts the back-end repositories that serve myriad tenant applications. For an analytics application, this would be the databases\/data warehouses\/data lakes that retain the underlying business data, together with any processed data objects for reporting and analysis. These objects are typically accessed by tenant users through reporting dashboards and data visualisation suites, abstracting users away from the complexities of database organisation.\n\\begin{figure}[h]\n \\includegraphics[width=\\linewidth]{.\/figures\/Cloud_layers_modeled_in_this_research.jpg}\n \\caption{Cloud layers modelled in this work.}\n \\label{fig:cloudlay}\n\\end{figure}\n\\subsection{Simulation design}\nWe have elected to model and simulate the proposed system (Figure \\ref{fig:cloudlay}) so that we can examine the operational characteristics in terms of performance, and its ability to provide a collection of services that are resilient towards malicious attacks.\n\n\n\nFigure \\ref{fig:virmac} illustrates how the cloud layers have been mapped to individual profiles. Each profile consists of a collection of VMs, that host the contents of the model as described in the previous section, namely: security and privacy services, tenant applications and repositories.\n\\begin{figure}[h]\n\\centering\n \\includegraphics[width=0.85\\linewidth]{.\/figures\/Virtual_Machines_with_applications_packaged.jpg}\n\\caption{Virtual Machines with applications packaged.}\n \\label{fig:virmac}\n\\end{figure}\n\\section{Modelling adversarial attacks}\nIf we now consider an adversarial attack upon the model, we can see how the system protects against such a scenario.\nFigure \\ref{fig:attack} shows the situation where a malicious agent attempts to infiltrate the system to obtain access to an authorised tenant's VM. Whilst the cloud service requires verification to be able to subscribe and access the remote resources, what appears to be a legitimate tenant could actually be an adversary that is masquerading as a valid client, who has the objective of entering the cloud and then attacking other tenant VMs from within the same cloud.\n\nTools such as Metasploit can be employed to automate the delivery of exploits in a rapid fashion. This could enable a bogus tenant to create surreptitious means of exposing sensitive data, unbeknown to any other party.\n\n\\begin{figure}[tb]\n\\centering\n \\includegraphics[width=0.9\\linewidth]{.\/figures\/An_attack_scenario_showing_the_problem.jpg}\n \\caption{An attack scenario showing the problem.}\n \\label{fig:attack}\n\\end{figure}\n\nPosing as a legitimate tenant, the adversarial agent would in this case be attacking from a VM that is authorised and hosted within the cloud. As such, the conventional cloud security processes would not be able to detect such activity.\nIn effect, the activity is obscured by the sheer volume of VMs that exist within a cloud environment.\nThis is a significant challenge for cloud service providers, particularly as service orientation through Microservice Architectures becomes more prevalent \\cite{hill2017}. If we extend this to the IoT domain, there is a stronger desire to package functionality into services, to be hosted on distributed, connected hardware. Therefore, the ability to successfully address this issue is a key feature of this work.\nWe propose a solution whereby the collection of VMs are organised into a hierarchy, as Figure \\ref{fig:ahier} shows.\nIf an adversarial agent enters the cloud via a subscription, and is assigned VM2, it has the potential to employ cross-channel attacks against VM1 and VM3, thereby exploiting the presence of virtual links between differing VMs. Typically, cloud security controls are deployed to prevent external attacks rather than insider attacking.\n\\begin{figure}[tb]\n \\includegraphics[width=\\linewidth]{.\/figures\/A_hierarchical_framework_presented_as_a_solution_to_the_attack_scenario_in_Figure_4.jpg}\n \\caption{A hierarchical framework presented as a solution to the attack scenario in Figure \\ref{fig:attack}.}\n \\label{fig:ahier}\n\\end{figure}\nThis architecture limits the opportunity to commit further exploits since the attacker is prevented from moving to the next control using virtual links. The only option is to proceed using a real network link by requesting a session in Control A, to communicate with VMs 4,5 and 6. Since there now exists Control A, the attacker has to successfully satisfy {\\tt Tenant Metadata Inspection} in order to proceed with the penetration.\nOf course, an orchestrated and sustained attempt to commit an attack will mean that we should anticipate an adversarial agent will also have obtained valid credentials, either by posing as a legitimate tenant or otherwise.\nIn this case, Control B will have required that the malicious agent would need to navigate the entire stack of cloud layers before access could be gained to the analytics interface.\n\nIf the attacker has satisfied the cloud validation and metadata inspection layers, the only way forward now is to plant exploits in the hope that these will lie undetected. However, the Intrusion Protection layer, and the Anti-Malware layer both offer protection for subversive, covert attacks from the inside.\nThe result is that our proposed model prevents data breaches, even when adversarial attacks are launched from what appears to be genuine service subscribers.\nWe can see in Figure \\ref{fig:control} a sequence in which various security controls might be instantiated. The security policies of the host system (or systems) will inform the order in which controls are implemented, to suit the goals desired by the infrastructure provider. It is also evident how a tenant's session is routed through the various VMs in order to access the relevant analytics services.\n\\begin{figure}[tb\n \\includegraphics[width=\\linewidth]{.\/figures\/The_controls_positioned_on_each_layer_of_the_hierarchy.jpg}\n \\caption{The position of the controls at each layer of the hierarchy.}\n \\label{fig:control}\n\\end{figure}\n\\section{Session packet inspection}\nIn this section we shall describe a detailed walkthrough of the security model and explain the use of session packet inspection.\n\nReferring back to Figure \\ref{fig:ahier}, VM1 will host a client that ultimately will access the enterprise analytics application hosted in VM7. A malicious agent that has access to either VM2 or VM3 (or both) can only attack the client of the analytics application, rather than the application itself.\nVM2 will be used to test the validity of the VM1 client using a VM identification number, and assuming that all is well, will launch a form via VM3 to request details from the tenant.\nThe details requested from the tenant will vary each time that VM3 is executed, but they will always refer to some personal details that can be used to help identify the correct tenant.\nOnce the form on VM3 has been completed, VM2 will use the responses to verify the details against those held in the {\\tt MetaDB} repository. Once the metadata has been verified, the session can continue. If the malicious tenant has the intention of acquiring metadata about other tenants, it would need to deliver an exploit into {\\tt MetaDB}. Whilst it might be expected that the system is fully patched, it is still conceivable that a vulnerability exists, enabling an adversary to progress to the next cloud layer.\nSince a session cannot be interrupted, any evidence left by malicious exploits will still remain as the agent will not be able to remove the incriminating evidence of the exploit. As such, the attacker can only proceed to the next layer by exposing that an exploit has been used to obtain entry. If the malicious agent has obtained private keys, they still cannot progress without exposing details of the exploit within the session.\nLayers 5 and 6 both have the capability to detect malware, which of course is reliant upon adequate, proactive security maintenance to ensure that all exploit databases are current.\n\nIt is feasible that an attacker could compromise VMs 3 and 4 with a fresh exploit that has yet to be discovered and documented, in which case the anti-malware cloud layer will be unaware of this exploit as well. Whilst this situation may foil malware detection, there still exists Intrusion Prevention within the cloud layer, which by its nature monitors and reports upon anomalies.\n\nOur model enables system architects the ability to add or subtract security controls as required for a given set of policies, merely by specifying additional cloud layers for the hierarchy. When a session is authenticated as satisfying the requirements of each layer, it can progress to subsequent layers until the destination application layer is reached.\nOne advantage of the use of multiple VMs is that much of the computation can be performed in parallel, and therefore the majority of packet inspection incurs a minimal overhead in VMs 2,3 and 4. However, session packet inspection across cloud layers, especially anything that involves Intrusion Prevention or anti-malware detection will result in an additional overhead. \n\\subsection{Mapping controls to the seven layer model}\nWith reference back to the NIST seven layer model \\cite{nist2011}, the controls of our proposed model can be mapped to the PaaS layer as per Figure \\ref{fig:mapp}, with the exception of the firewalls which naturally reside in the IaaS layer \\cite{hill2013,carvakho2011}.\nTenant VMs are hosted in layers 1-3 of the NIST model; users interact directly with clients that do not store data.\n\nSessions commence within tenant VMs, before moving to the application in layer 7 via several layers of verification and authentication in layers 4 and 5. If the tenant has a SaaS platform, this will be made available through layer 6, otherwise all bespoke applications are resident in layer 7.\n\\subsection{Session flow}\nThe proposed architecture is located on layers 4 (IaaS) and 5 (PaaS) of the NIST seven-layer model. All firewalls are categorised as IaaS due to the functionality of verifying VM instance IDs, and also relating these checks to the authorisation data provided by a tenant to secure access to the cloud.\nWhilst VM IDs are assigned in layers 2 and 3, access controls are assigned at layer 4. Subsequent controls are assigned outside of the VM layer, and are primarily concerned with session packet inspection, for a given session $S$.\nInformation is requested directly from the tenant to satisfy the DB\\textsubscript{META} and DB\\textsubscript{VAULT} checkpoint controls, and is supplemented by DB\\textsubscript{IPS} and DB\\textsubscript{ANTIMAL} controls that perform the session packet inspection function. DB\\textsubscript{IPS} and DB\\textsubscript{META} are therefore categorised as PaaS controls.\nIt is likely that controls will also exist at the application layer. For instance, a SaaS instance will require user authentication, as will a custom enterprise application\\cite{ouf2011}. User role profiles are useful in such scenarios to manage different levels of system access within an enterprise application, to reflect the role, responsibilities and authority of a particular stakeholder.\n\n\n\n\\begin{figure*}[htb]\n\\centering\n\\includegraphics[width=0.9\\textwidth]{.\/figures\/mapping.jpg}\n \\caption{Mapping of the proposed architecture with the NIST seven layer cloud model\\cite{nist2011}.}\n \\label{fig:mapp}\n\\end{figure*}\n\\section{Algorithm design}\nAlgorithm 1 describes the sequence of inspections that are performed within the proposed multi-layer security model. This algorithm represents the core functionality and can be augmented as additional cloud layers are defined in response to the security needs of an organisation. In the case of a single enterprise the model may tend towards fewer augmentations.\nHowever, enterprises that collaborate, or who choose to adopt shared services across distributed platforms, will no doubt adopt additional layers in order to securely manage service access. \n\\begin{algorithm}\n\\SetAlgoLined\nInput: $S$, $P_{CT}$, ($DB_{FW}$, $DB_{META}$, $DB_{VAULT}$, $DB_{IPS}$, $DB_{ANTIMAL}$), 1=permit, 0=deny\\\\\nTenant session: $S$\\\\\nContents of session packets:$P_{CT}$\\\\\nContents of FW: $DB_{FW}$\\\\\nContents of $TENANT_{META}: DB_{META}$\\\\\nContents of $TENANT_{VAULT}: DB_{VAULT}$\\\\\nContents of $IPS: DB_{IPS}$\\\\\nContents of $ANTIMALWARE: DB_{ANTIMAL}$\\\\\nFlags: 1=permit, 0=deny\\\\\n \n Initialise $S$\\;\n Set $S=1$, Match$(P_{CT})$;\\\\\n \\ForEach{\\{$DB_{FW}$, $DB_{META}$, $DB_{VAULT}$, $DB_{IPS}$, $DB_{ANTIMAL}$\\}} {%\n \\eIf{$P_{CT}\\in$ \\{$DB_{FW}$,$DB_{META}$, $DB_{VAULT}$\\} AND $P_{CT}\\notin$ \\{$DB_{IPS}$,$DB_{ANTIMAL}$\\}}{\n set $S=1$; AuthoriseTenantAccess()\\; \/\/tenant access authorised\\;}{\n set $S=0$; DenyTenantAccess()\\; \/\/tenant access denied\\;\n }\n}\n Output: $S$\n \\caption{Multi-layer hierarchical packet inspection}\n\\end{algorithm}\n\\subsection{Security model logic}\nOur proposed model seeks to address security concerns in distributed service applications across heterogeneous hardware resources by enabling session packet inspection to take place at a number of checkpoints. Each session packet is scrutinised and compared with a number of repositories such as DBMETA, DBVAULT, etc. Each inspection stage shall now be described in turn.\n\\begin{enumerate}\n\\item {\\bf DB\\textsubscript{FW}:} For each instance, a client initiates a session, which is inspected as a second stage. This session will possess a VM ID together with authentication and authorisation verification data.\nIn addition, the DB\\textsubscript{FW} must have an entry that relates to the VM ID, otherwise the session is terminated.\n\\item {\\bf DB\\textsubscript{META}:} A further inspection is then performed to confirm tenant metadata as requested by the cloud host administrators.\nOnce the metadata has been provided, this will be added to session packets. The session can only continue if the session metadata matches that which exists in the DB\\textsubscript{META} repository.\n\\item {\\bf DB\\textsubscript{VAULT}:} At this stage the VM ID within the session will be inspected to verify that a private key exists within the DB\\textsubscript{VAULT} so that encrypted databases can be accessed.\n\\item {\\bf DB\\textsubscript{IPS}:} On reaching this stage, the session itself is now regarded as being fully authenticated and is authorised to progress to subsequent layers. The next stage is to inspect the sessions for evidence of potential exploits that match those held in the DB\\textsubscript{IPS} repository. If a match occurs, the session is terminated.\n\\item {\\bf DB\\textsubscript{ANTIMAL}:} This stage performs an anti-malware check, which in conjunction with the DB\\textsubscript{IPS} inspection, prevents fully authenticated VM sessions being able to penetrate the upper layers of the security model by masquerading as legitimate cloud tenants. \n\\end{enumerate}\nThe algorithm ensures that after a session is initiated, it is inspected at each layer, where each layer is represented by a separate cloud. A key principle is that session packet data must match the firewall data, tenant metadata and vault data before a session can be authorised. Furthermore, the session cannot be granted access to application layers until it has been successfully screened against IPS and anti-malware repositories.\n\\subsection{Implementing the model}\nFor our example scenario, the LAN contains 500 clients. Each of the clients is assigned three VMs, with an assigned destination being the tenant client's metadata repositories rather than the eventual analytics application. This was a conscious decision to prohibit any tenant sessions from attempting to subvert the security and privacy controls of the cloud.\nFor a physical network, this control would most likely be enacted by a firewall in an application layer, or as part of a setting in a virtual network controller. For consistency we have replicated this by ensuring that the destinations of the metadata servers are directed towards the tenant vaults; this is a faithful representation of the intentions of the algorithm, in that it is governing the marshalling of each session packet by enforcing checkpoints in a particular sequence. In the case of a session packet not containing a tenant key that matches a corresponding entry in DB\\textsubscript{VAULT}, the encryption key is not assigned and the packets are dropped \\cite{diaz2013}. We have represented the adversarial agents as clients who have the authorisation of a valid tenant, in order to simulate an attack from within.\n\\section{Results}\nFigure \\ref{fig:dbsess} shows the sessions that have been hosted by the tenant LAN in the simulation. We observed that sessions only existed within the tenant LAN for tenants that had matching metadata in DB\\textsubscript{META}. It follows that the sessions that were invoked within the tenant LAN, had consistency between tenant metadata (DB\\textsubscript{META}), tenant vault (DB\\textsubscript{VAULT}), and the associated VM ID from the initial authentication.\nThe simulation demonstrated that these validations were maintained throughout the experiment for all network hops, illustrating that VMs were complying with their specified destination configurations, and there was no evidence of VMs by-passing authentication layers. As such, all security and privacy control rules are enforced by the model.\n\n\\begin{figure}[tb\n\\centering\n \\includegraphics[width=0.9\n \\linewidth]{.\/figures\/Client_DB_sessions_on_Tenants__LAN.jpg}\n\\caption{Client DB sessions on Tenants' LAN.}\n \\label{fig:dbsess}\n\\end{figure}\nVMs initiated by malicious agents have separate profiles from those of legitimate tenants, even though a malicious agent may appear legitimate at the outset. This means that adversaries had distinct metadata and decryption keys in their vault repositories.\nWe can see in Figure \\ref{fig:ipdropped} the instances where a malicious agent's packets have been dropped as a result of packet inspection through either the IPS or anti-malware layers, preventing deeper penetration into the system. As discussed earlier, an internal attack (an agent with legitimate DB\\textsubscript{META} data) would need to compromise DB\\textsubscript{VAULT}, DB\\textsubscript{IPS} and DB\\textsubscript{ANTIMAL} in sequence if it is to successfully reach the analytics application layer.\n\\begin{figure}[tb\n \\includegraphics[width=\\linewidth]{.\/figures\/IP_packets_from_the_hackers__machines_are_dropped.jpg}\n \\caption{After initial attempts to penetrate the system, IP packets from the attackers' machines are dropped.}\n \\label{fig:ipdropped}\n\\end{figure}\nConversely, an authorised tenant may elect to execute DB sessions on their own LAN, since the application profiles can include references to their VMs, as per Figure \\ref{fig:auth}. In such cases, the sessions are authorised in the sense that they fulfil the relevant rule in Algorithm 1. Since we have chosen to model each of the layers as clouds, each of the multiple layers can in fact be serviced by different cloud providers. This architecture thus demonstrates significant flexibility and is attractive to system architects who are proactively designing systems that will rely upon heterogeneous hardware and distributed resources, such as the IoT and IIoT environments.\nAn enterprise may adopt this security model so that it can take the opportunity to employ software applications and services that are themselves hosted on offsite private or community clouds. However, the use of these services may increase operational costs, although this is off-set by a reduction in capital expenditure. Operational costs may also increase indirectly through the maintenance charges associated with managing the database updates of five security control layers.\n\\begin{figure}[tb\n \\includegraphics[width=\\linewidth]{.\/figures\/Authorized_tenant_LANs_established_and_ran_DB_sessions_with_the_TENANT_METADATA}\n \\caption{Authorised tenant LANs established, DB sessions initiated with the TENANT METADATA.}\n \\label{fig:auth}\n\\end{figure}\n\\section {Conclusions\nThis article proposes a multi-layer hierarchical inter-cloud security model, that is inspired by the NIST seven-layer model of cloud computing. By using sequential session packet inspection techniques, we have demonstrated an architecture that exhibits considerable resilience towards both external attack as well as more surreptitious internal adversarial behaviour. Whilst VM vulnerabilities are well documented in multi-tenant shared environments, our five layers for packet inspection enables the architecture to identify and compartmentalise malicious activity. Thus, penetration of a firewall is in itself insufficient as a means of attempting to access the application layer, as it is then necessary to create an evidence trail of exploits that cannot be hidden from the IPS and anti-malware packet inspection layers. The proposed model is particularly suited to architectures that have a requirement to remain flexible for future scaling (which is often a driver for the adoption of cloud infrastructure), such as those built upon microservices. \nOur solution is mapped to the NIST model in order to assist cloud and IoT system architects to incorporate this work into their own designs.\n \n \n \n \n \n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nCompact groups of galaxies are associations of three to seven galaxies, where the projected distances between them is of the order of their diameters, and where the group\nshows a low velocity dispersion, making compact groups an ideal place to study\ngalaxy interaction and intergalactic star formation (e.g. \\cite[Torres-Flores et\nal. 2010]{tor09b}, \\cite[de Mello et al. 2008]{dem08b}, \\cite[de Mello, Torres--Flores \\& Mendes de Oliveira 2008]{dem08a}, \\cite[Mendes de Oliveira et al. 2004]{mdeo04}). The main goal of this work is to search for a link between the evolutionary stage of a group and the presence of young intergalactic objects which may have formed during galaxy interactions. For this, we analyze a subsample of seven compact groups (HCG 2, 7, 22, 23, 92, 100 and NGC 92) which span a wide range of evolutionary stages, from HI rich groups to strongly interacting groups, where the galaxies show tidal tail features and a deficiency in neutral HI gas. In order to analyze the evolutionary stage of each group, we used new Fabry-Perot velocity maps, GALEX\/UV data and optical R-band images. The velocity fields and rotation curves help constraining the evolutionary stage of each compact group while ultraviolet light contains important information regarding the age of the young stellar population that may be present in the intragroup medium.\n\n\\section{UV analysis}\n\nWe searched for ultraviolet emitting regions in the vicinity of all seven targets, using the SExtractor software (SE, \\cite[Bertin \\& Arnouts 1996]{ber96}) in the FUV, NUV and R sky-subtracted images of our compact group sample. We compare the field density of regions detected in the compact group with a control sample outside the group. HCG 92 and HCG 22 have the highest field density in this study. No excess was found in HCG 2, HCG 7, HCG 23, HCG 100 and NGC 92 (\\cite[Torres-Flores et al. 2009]{tor09}).\n\n\\section{Fabry-Perot analysis}\n\nIn order to constrain the evolutionary stage of each compact group, we\ninspected the velocity field and rotation curve of each galaxy to search for\ninteraction indicators, in a similar way to that done by \\cite[Plana et al. (2003)]{pla03} and \\cite[Amram et al. (2003)]{amr03}. In the case of NGC 92, it shows a prominent tidal tail in its velocity field. At the tip of this tail, there is a tidal dwarf galaxy candidate, having an age of about 40 Myrs (\\cite[Torres-Flores et al. 2009]{tor09}).\n\n\\section{Conclusions}\n\nWe used multiwavelength data to study the evolutionary stages of the compact groups of\ngalaxies HCG 2, 7, 22, 23, 92, 100 and NGC 92. New Fabry-Perot velocity fields, rotation curves and GALEX NUV\/FUV images were analyzed for four and seven of these groups respectively. Groups HCG 7 and 23 are in an early stage of interaction whereas\nHCG 2 and 22 show limited interaction features and HCG 92, 100 and NGC 92 are in a\nlate stage of evolution, having HI gas in the intragroup medium, galaxies with peculiar velocity fields and several young star-forming regions in the intergalactic medium.\n\n\\acknowledgments\n\nS. T--F. acknowledges the financial support of FAPESP through the Doctoral position, under contract 2007\/07973-3 and Eiffel scholarship.\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzbjtv b/data_all_eng_slimpj/shuffled/split2/finalzzbjtv new file mode 100644 index 0000000000000000000000000000000000000000..324d79278c765271d3381e17eb0af940807e4ffc --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzbjtv @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\nChemical signalling that exploits molecule or ion for communication has been widely found in nature with diverse propagation distances. In a relatively long range, animals or insects utilize pheromone to communicate with the members of their species for mate selection, identity recognition, alarm inform, etc. \\cite{Shorey.Pheromone.13,Akyildiz13}. At micro-scale environment, hormones or other chemical substances are prevalently transmitted or received in tiny organisms such as cells. This process is the so-called cell signalling that is crucial to cells' survival~\\cite{Han.Signalling.17}.\n\nMolecular communication (MC) is an emerging technique inspired from the aforementioned communication schemes in vivo, whose history dates back to 2005 \\cite{Tadashi.MC.05}. MC is available at both macro-scale and micro-scale~\\cite{Farsad.MCSurvey.16}. At micro-scale, MC is suitable for connection among nanomachines, whose communication distance ranges from a nanometer to a hundred nanometers. Nanomachine is one of the most remarkable progress of nanotechnology that has the potential to revolutionize many aspects of technology, and ultimately benefit our life. However, a single nanomachine can only perform very simple tasks due to its size constraint, whereas complex applications including biopsy, targeted drug delivery, environmental sensing, cell sorting, etc., require coordination among a swarm of nanomachines~\\cite{Wang.Nanomachine.13}. MC is regarded as a prominent candidate for nanonetworking because of its bio-compatible and energy-efficient characteristics. It has been acknowledged as the most important communication paradigm in IEEE 1906.1 standard~\\cite{Tadashi.10year.17}.\n\nModulation plays a significant role in MC as it determines the system's achievable performance~\\cite{Suzuki.MCmodelling.17}. \nThere are a few of modulation schemes for MC proposed in the literature, including the concentration shift keying (CSK)~\\cite{Kuran.ModuMC.11}, molecular shift keying (MoSK)~\\cite{Kuran.ModuMC.11}, isomer-based ratio shift Keying (IRSK)~\\cite{Kim.Isomodu.13} and pulse position modulation (PPM)~\\cite{Llatser.PPM.13}, in which messages are encoded as the concentration, type, ratio and release time of transmitted molecules, respectively. Advanced modulation schemes have also been considered in the context of the single-input single-output~(SISO) based MC (SISO-MC)~\\cite{Arjmandi.MCSK.13,Kabir.DMoSK.15,Mosayebi.typesignmodu.18}. However, SISO-MC is not always suitable for the scenarios of high-speed transmission and some of the applications where reliability is indispensable. To solve these problems, some well-known techniques in the conventional wireless communications, e.g., multiple-input multiple-output (MIMO), have been redesigned for MC \\cite{KC.frontierWC.12}. Specifically, MIMO based MC (MIMO-MC) is a recent trend in MC research, dated back to 2012 when it was first proposed in \\cite{KC.MIMOMC.12}. It is shown that in MIMO-MC, spatial diversity enhances the bit error rate (BER) performance, while spatial multiplexing may increase transmission rate significantly~\\cite{KC.MIMOMC.12}. In 2013, a micro-scale MIMO-MC system was introduced~\\cite{Tadashi.MCbook.13}, where a group of sender nanomachines simultaneously transmit messages to a group of receiver nanomachines through the medium where they reside. In the same year, the first SISO-MC prototype at macro-scale was implemented \\cite{Farsad.Prototype.13}, which mentioned that MIMO principle may be introduced to improve the transmission rate. In 2016, a $2 \\times 2$ MIMO-MC prototype achieving spatial multiplexing was implemented, which achieves a 1.7 times higher data rate than its SISO counterpart~\\cite{Chae.MCMIMOproto.16}. These research and practice demonstrate that spatial multiplexing in MIMO-MC is feasible for rate increase, although the data rate is not doubled due to the existence of interference and overhead. In~2017, MIMO-MC technique gained more attention than ever before, due to the appearance of the training-based channel estimation~\\cite{Rou.CEMIMOMC.17}, spatial diversity coding techniques~\\cite{Damrath.Spatialcoding.17} and the introduction of machine learning based channel modeling methods~\\cite{Lee.MLMIMOMC.17}. Moreover, in \\cite{WG.MCSIMOsyn.18}, synchronization was investigating in the context of the single-input multiple-output (SIMO) based MC (SIMO-MC). Now, we have no doubt that MIMO-MC constitutes a promising technique for performance improvement in MC. However, a typical challenge in MIMO-MC is the inter-link interference (ILI) in addition to the inter-symbol interference (ISI) existing in all MC systems. Due to the ILI, first, it increases the detection complexity at the receiver side. Second, the study in \\cite{Chae.MCMIMOproto.16}~manifests that the BER of $2 \\times 2$ MIMO-MC prototype is 2.2 times higher than that of its SISO-MC counterpart, meaning that high-reliability MC is challenging. Third, the energy consumption in MIMO-MC is significant in comparison to SISO-MC, due to more devices being simultaneously activated at both transmitter and receiver sides. This problem can be serious in micro-scale MC, since the power supply of nanomachine is limited and their computing capability is low. Consequently, the development of energy-efficient and low-complexity transmission schemes for MIMO-MC is demanding.\n\nIn this paper, we propose a spatial modulation based MC (SM-MC) for MIMO-MC implementation, which combines a space-dependent modulation with a concentration-relied modulation. By allowing only the space modulation, we also propose the special case of SM-MC, namely, the space shift keying based MC (SSK-MC). It can be shown that both SM-MC and SSK-MC are able to combat the aforementioned problems of the MIMO-MC. It is well known that SM~\\cite{Mesleh.SM.08} and SSK~\\cite{Jeganathan.SSK.09} in radio-based wireless communications exhibit low complexity and high energy efficiency~\\cite{Wen.IMsurvey.17}. By exploiting the spatial domain for message encoding via activating an active transmit antenna in each time slot, SM and SSK are capable of eliminating the inter-antenna interference (IAI). Note that, the major distinction between SM and SSK is that SM uses both spatial and signal constellation to transmit information, while SSK only exploits the spatial constellation.\n\nIn our design for SM-MC and SSK-MC, the spatial domain is reflected by the concentration difference at the receiver side, when different transmitters emit molecules. By contrast, the signal constellation domain is contributed by the different constellation level in CSK. Hence, the SM-MC\/SSK-MC receivers can collaborate to identify the activated transmitter nanomachine for detecting the space symbol, while the SM-MC receiver can decode the concentration symbol according to the rule of CSK. Furthermore, in order to enhance the reliability of detection, detector assisted by either \\emph{equal gain combining} (EGC) or \\emph{selection combining} (SC) is developed for detecting the concentration symbol. In this paper, we derive the SER expressions for both the SM-MC and SSK-MC, and conduct computer simulations to examine their performance as well as to validate our theoretical analysis. Our studies and results show that both the SM-MC and SSK-MC outperform the MIMO-MC in terms of the SER performance, and our analytical results are accurate enough for predicting the achieving SER performance of the SM-MC and SSK-MC systems.\n \nThe remainder of this paper is organized as follows. Section~II reviews the system model of MIMO-MC over diffusion channels. In Section~III, we present the principles of SM-MC and SSK-MC based on the architecture of MIMO-MC. In this section, we also derive the resulting SER of the SM-MC and SSK-MC. Section~IV evaluates the SER performance of SM-MC and SSK-MC by invoking the MIMO-MC with on-off keying (OOK) modulation and the SISO-MC with quadruple CSK (QCSK) modulation as the benchmarks. Finally, the research is concluded in Section~V.\n\n\\textit{Notation}: Boldface uppercase and lowercase letters indicate matrices and vectors, respectively. $\\mathbb{R}^{n\\times m}$ indicates a real number matrix with $n\\times m$ dimensions. $\\mathbb{E}[\\cdot]$, $\\left \\lVert \\cdot \\right \\rVert$, $|\\cdot|$, $Q(\\cdot)$ and $\\text{Pr}[\\cdot]$ represent expectation, Euclidean norm, absolute value, Q-function and probability of an event, respectively.\n\n\\section{Review of MIMO-MC}\nWe consider a MIMO-MC system in a 3-D unbounded environment with point sources and spherical receivers, which are assumed to be memoryless as \\cite{Yilmaz.14,Zamiri.CDMAMC.16,Tepekule.Memoryless.15}, meaning that they have no information of the previously detected symbols. Perfect synchronization is also assumed in our MIMO-MC. In this section, we propose the channel and communication models for MIMO-MC, which constitutes the fundamentals of our subsequent SM-MC scheme and its special case of the SSK-MC scheme.\n\\subsection{Channel Model of MIMO-MC}\n\\begin{figure}[tb]\n\t\\centering\n\t\\includegraphics[width=3.25in]{Fig1MIMOMC.eps}\n\t\\caption{System model of MIMO-MC.}\n\t\\label{fig_1}\n\\end{figure}\n\nWe consider an $N \\times N$ diffusion based MIMO-MC system operated at micro-scale in this paper \\cite{Tadashi.MCbook.13}, where $N$ transmit and receive nanomachines are attached to the cell membrane at both transmitter and receiver sides, respectively. Furthermore, we assume that there are a transmission control center and a reception decision center, respectively, located at the centers of the transmitter and receiver cells. The function of transmission control center is to coordinate the transmitter nanomachines to emit molecular pulses according to the information to be transmitted, while the reception decision center connecting all the receiver nanomachines decodes the information based on the received signals from $N$ receiver nanomachines. To clearly demonstrate the communication system model of the MIMO-MC, we exemplify a $4 \\times 4$ MIMO-MC system with $4$ pairs of transceivers as shown in Fig. 1, where the lipid bilayer of cell membrane are shown in green color. As shown in Fig. 1, the pair of transmitter and receiver of a link and the molecules transmitted between them are marked using the same color for the sake of easy distinction. Under the assumption that perfect alignment is achieved in the MIMO-MC architecture, the spacing between adjacent transmitters or adjacent receivers is equally set as $r$, and the distance between a pair transmitter and receiver is expressed as $d$. Thus, the distance from the $j$-th transmitter to the $i$-th receiver is given by\n\\begin{align}\nd_{ji} =\n\\begin{cases}\n\\hspace{1cm}d, & \\text{for $i=j$}, \\\\\n\\sqrt{d^{2}+|j-i|^{2}r^{2}}, & \\text{for $i\\neq j$.}\n\\end{cases}\n\\end{align}\n\nIn this paper, we assume $M$-ary CSK modulation and denote the concentration of molecules in the $i$-th receiver at time $t$ in response to the $j$-th transmitter as $c_{m,ij}(t)$ when a pulse of $S_{m}$ molecules with $m\\in\\{0,1,\\ldots, M-1\\}$ is emitted at $t=0$. Consequently, the concentration at receiver can be formulated according to Fick's second law of diffusion as \\cite{Dk.receiverMC.13}\n\\begin{align}\nc_{m,ij}(t)=S_{m}\\frac{1}{{(4\\pi Dt)}^\\frac{3}{2}}\\exp \\left(-\\frac{d_{ji}^2}{4Dt}\\right),\n\\begin{array}{c}\n\\hspace{-0.3cm}j\\in\\{1,2,\\ldots, N\\},\\\\\n\\hspace{-0.3cm}i\\in\\{1,2,\\ldots, N\\},\\\\\n\\hspace{-0.4cm}m\\in\\{0,1,\\ldots, M-1\\},\\\\\n\\end{array}\n\\end{align}\nwhere $D$ is the diffusion coefficient of the propagation medium that is assumed to be homogeneous in this paper. We assume that the spherical receiver nanomachine has a volume of $V_{\\sss \\text{RX}}=\\frac{4}{3}\\pi{\\rho}^{3}$ with $\\rho$ being the radius of the receiver. Molecular concentration is assumed to be uniform inside a passive receiver when $d\\gg\\rho$. Based on these assumptions, the expected number of molecules inside the $i$-th receiver at time $t$ can be formulated by \\cite{MM.studyofCEMC.14,ANoel.Passivereceiver.16}\n\\begin{align}\nN_{m,ij}(t)=V_{\\sss\\rm{RX}}c_{m,ij}(t)=V_{\\sss\\rm{RX}}S_{m}h_{ij}(t),&\\text{\\quad\\quad for $t>0,$}\n\\end{align}\nwhere $h_{ij}(t)$ indicates the probability that an information molecule released at $t=0$ is sensed by the passive receiver at time~$t$ \\cite{GeC.adadetec.18,GD.AmpliDetect.17}, which is given by\n\\begin{align}\nh_{ij}(t)=\\frac{1}{{(4\\pi Dt)}^\\frac{3}{2}}\\exp\\left(-\\frac{d_{ji}^2}{4Dt}\\right).\n\\end{align}\nIn MC, $h_{ij}(t)$ represents the channel impulse response (CIR) between the $j$-th transmitter and the $i$-th receiver. Therefore, the concentration vector $\\mathbf{c}_{m,j}(t)$$\\in \\mathbb{R}^{N \\times 1}$ collecting the expected concentration of all receivers in response to the $j$-th transmitter can be written as\n\\begin{align}\n\\mathbf{c}_{m,j}(t)&\\triangleq\\big[c_{m,1j}(t), \\ldots, c_{m,ij}(t), \\ldots, c_{m,Nj}(t)\\big]^{T}\\nonumber\\\\\n&=S_{m}\\big[h_{1j}(t), \\ldots, h_{ij}(t), \\ldots, h_{Nj}(t)\\big]^{T}\\nonumber\\\\\n&=S_{m}\\mathbf{h}_{j}(t),\n\\end{align}\nwhere $\\mathbf{h}_{j}(t)$ is the CIR vector from the $j$-th transmitter to all the $N$ receivers at time $t$. Let us define $\\mathbf{H}(t)$$\\in \\mathbb{R}^{N \\times N}$ as the channel matrix of an $N \\times N$ MIMO-MC system at time $t$, which can be represented as\n\\begin{align}\n\\mathbf{H}(t)&=\\Big[\\mathbf{h}_{1}(t); \\ldots ;\\mathbf{h}_{j}(t); \\ldots ; \\mathbf{h}_{N}(t)\\Big],\n\\end{align}\nwhose entries are given by (4). Based on our assumptions, we can know that the diagonal elements of $\\mathbf{H}(t)$, such as $h_{jj}(t)$, have the same value, given by (4) associated with setting $i=j$ and $d_{jj}=d$. By contrast, a non-diagonal element $h_{ij}(t)$ gives the probability of the ILI from the transmitter $j$ to the receiver~$i$.\n\nIn this paper, we consider the amplitude detection \\cite{Llatser.PPM.13}. To achieve this, we assume that all the receivers sense the concentration at a certain time, e.g., at the time when the concentration at a receiver generated by its paired transmitter reaches the peak value, which can be obtained by solving the partial derivative equation $\\frac{\\partial{c_{m,jj}(t)}}{\\partial{t}}=0$. Specifically, when an impulse of molecules is emitted by a transmitter at $t=0$, the peak concentration presenting at its paired receiver can be found to be at the time of\n\\begin{align}\nt_{p}=\\frac{d^{2}}{6D}.\n\\end{align}\nExplicitly, the peak time is irrelevant to $S_{m}$. In this paper, without further explanation, it is assumed that each receiver of the MIMO-MC samples for concentration after a time interval $t_{p}$ seconds from the emission of the chemical impulse by its paired transmitter. Therefore, upon substituting (1) and (7) into (2), we can derive the maximum concentration of a receiver in response to its paired transmitter as\n\\begin{align}\nc_{m,jj}(t_{p})=S_{m}\\bigg(\\frac{3}{2\\pi ed^{2}}\\bigg)^{\\frac{3}{2}},\n\\end{align}\nwhen an impulse of $S_{m}$ molecules is released. Similarly, $c_{m,ij}(t_{p})$ can be derived, given by\n\\begin{align}\nc_{m,ij}(t_{p})=S_{m}\\bigg(\\frac{3}{2\\pi d^{2}}\\bigg)^{\\frac{3}{2}}\\exp \\left(-\\frac{3d_{ji}^2}{2d^2}\\right).\n\\end{align}\n\n\\subsection{Communication Model of MIMO-MC}\n\\begin{figure}[!t]\n\t\\centering\n\t\\includegraphics[width=3.25in]{Fig2ILIISIinMIMO_MC.eps}\n\t\\caption{Concentration expected at one receiver of a $2\\times2$ MIMO-MC affected by its ISI and the ILI from one other link, when $\\text{SNR}=10\\ \\rm{dB}$, $D=2.2\\times10^{-9}$$m^{2}\/s$, $d=20um$, $r=15um$ and $T_{s}=0.2s$.}\n\t\\label{fig_2}\n\\end{figure}\n\nBinary CSK (BCSK) is the simplest CSK with $M=2$, which emits a chemical pulse containing $S_{1}$ molecules towards its paired receiver for transmitting bit ``1'', or a pulse of $S_{0}$ molecules for transmitting bit ``0''. As a special case of BCSK, the OOK modulation, which is prevalently adopted in the existing MIMO-MC \\cite{KC.MIMOMC.12,Chae.MCMIMOproto.16,Lee.MLMIMOMC.17,Damrath.Spatialcoding.17}, keeps silent without any molecule emission for bit ``0'', i.e., $S_{0}=0$. For the general CSK, it has been revealed in literature that the CSK with $M\\geq4$ is usually unable to attain satisfactory error performance in SISO-MC \\cite{Suzuki.MCmodelling.17,Kabir.DMoSK.15}. Note that in Fig. 1, we considered OOK as the modulation scheme for MIMO-MC, where the red, yellow, and purple colored links represent transmitting $S_{1}$ molecules to their corresponding receivers, while the blue colored one keeps silent without any emission of molecules. Furthermore, in Fig. 1, there are still a few of blue-colored molecules, which are the residual molecules of the previous transmissions. \n\nLet us denote the transmit signal vector at the sampling time $t$ as\n\\begin{align}\n\\mathbf{x}(t)=\\Big[x_{1}(t),\\ldots,x_{i}(t),\\ldots, x_{N}(t)\\Big]^{T},\n\\end{align}\nwhere $x_{i}(t)$$\\in \\{S_{0},S_{1}\\}$ denotes the number of molecules emitted by the $i$-th transmitter. Then, based on (6) and (9), the concentration vector corrupted by noise sensed at time $t$ can be expressed as\n\\begin{align}\n\\mathbf{y}_{\\sss \\text{MIMO}}(t)&=\\Big[{y}_{1,\\sss \\text{ MIMO}}(t),\\ldots,{y}_{i,\\sss \\text{ MIMO}}(t),\\ldots, {y}_{N,\\sss \\text{ MIMO}}(t)\\Big]^{T}\\nonumber\\\\\n&=\\sum_{l=0}^{L}\\mathbf{H}(t+lT_{s})\\mathbf{x}(t-lT_{s})+{\\mathbf{n}_{\\sss \\text{MIMO}}}(t),\n\\end{align}\nwhere $y_{i,\\sss \\text{ MIMO}}(t)$ represents the concentration sensed by the $i$-th receiver at time $t$, $T_{s}$ is the symbol duration. We assume that the ISI and ILI last for $L$ and $L+1$ symbol durations, respectively, where ISI is generated by the paired transmitter, while the ILI is the interference resulted from the other links. Practically, there is also noise in MC, which is given by the form of the noise vector ${\\mathbf{n}_{\\sss \\text{MIMO}}}(t)$$\\in \\mathbb{R}^{N \\times 1}$, having\n\\begin{align}\n\\mathbf{n}_{\\sss \\text{MIMO}}(t)=[n_{1,\\sss \\text{ MIMO}}(t),\\ldots,n_{i,\\sss \\text{ MIMO}}(t),\\ldots, n_{N,\\sss \\text{ MIMO}}(t)]^{T},\n\\end{align}\nwhere $n_{i,\\sss \\text{ MIMO}}(t)$ is a signal dependent noise sensed by the $i$-th receiver. Similar to \\cite{Damrath.Spatialcoding.17}, $y_{i,\\sss \\text{ MIMO}}(t)$ is the concentration sensed by the $i$-th receiver at time $t$:\n\\begin{align}\ny_{i,\\sss \\text{ MIMO}}(t)= \\sum_{j=1}^{N}\\sum_{l=0}^{L}{h}_{ij}(t+lT_{s}){x}_{j}(t-lT_{s})+n_{i,\\sss \\text{ MIMO}}(t).\n\\end{align}\n\nThe impact of ISI and ILI on a desired signal is illustrated in Fig. 2 for a $(2\\times2)$ MIMO-MC system, where the ILI is not only from the previous molecular symbols, but also from the current one released by the unpaired transmitter. However, it is worth noting that the current ILI plays a significant role in confusing the detection of a desired molecular signal, since the current ILI, as shown in Fig. 2, generates the strongest interference. The existing studies on SISO-MC often ignore the effect of ISI or only take one previous symbol into account \\cite{AN.Enzyme.14}, due to the fact that this previous symbol generates the highest ISI. However, when MIMO-MC is considered, as shown in Fig.~2, the effect of the current ILI may significantly surpass that of the ISI. In other words, the interference in MIMO-MC system may be dominated by the current ILI generated by the other unpaired transmitters, and the ISI generated by one previous symbol transmitted by the paired transmitter. Therefore, if we express (13) as\n\\begin{align}\ny_{i,\\sss \\text{ MIMO}}(t)=\n\\underbrace{{h}_{ii}(t){x}_{i}(t)}_{\\text{desired signal}}+\\underbrace{I_{i,\\sss \\text{ MIMO}}(t)}_{\\text{sum of interference}}+\\underbrace{n_{i,\\sss \\text{ MIMO}}(t)}_{\\text{noise}},\n\\end{align}\nthe noise component can be assumed to follow the Gaussian distribution, depending on the current symbol at time $t$ as \\cite{MengLS.receiverdesign.14}:\n\\begin{align}\nn_{i,\\sss \\text{ MIMO}}(t)\\sim \\mathcal{N}\\big(\\mu_{ni,\\sss \\text{ MIMO}}(t),\\sigma^{2}_{ni,\\sss \\text{ SM}}(t)\\big),\n\\end{align}\nwith\n\\begin{align}\n\\mu_{ni,\\sss \\text{ MIMO}}(t)=0,\\quad \\sigma^{2}_{ni,\\sss \\text{ SM}}(t)=\\frac{{h}_{ii}(t){x}_{i}(t)+I_{i,\\sss \\text{ MIMO}}(t)}{V_{\\sss \\text{RX}}}.\\nonumber\n\\end{align}\n\nThe interference component in (14) can be approximated as\n\\begin{align}\nI_{i,\\sss \\text{ MIMO}}(t)\\approx\\underbrace{\\sum\\nolimits_{j\\neq i}{h}_{ij}(t){x}_{j}(t)}_{\\text{current ILI}}+\\underbrace{{h}_{ii}(t+T_{s}){x}_{i}(t-T_{s})}_{\\text{last ISI}},\n\\end{align}\nwhen considering that the ILI is only from the current unpaired transmitter, while the ISI is only from one previous symbol sent by the paired transmitter.\n\\section{Spatial modulation based molecular communication}\n\nFrom the above description, we can know that MIMO-MC experiences both ILI and ISI, which may significantly degrade the detection performance. In order to combat these problems, in this section, we propose the spatial modulation based MC (SM-MC) as one of the implementation of the MIMO-MC. In our proposed scheme, both spatial and concentration domains are exploited for conveying information simultaneously, but only one type of molecules is used for transmission.\n\\subsection{Transmitter of SM-MC}\n\\begin{figure}[t]\n\t\\centering\n\t\\includegraphics[width=3.2in]{Fig3DiagramSM_MC.eps}\n\t\\caption{Transmitter diagram of SM-MC.}\n\t\\label{fig_3}\n\\end{figure}\n\nThe ideology of SM-MC is inspired by the SSK technique having been widely studied in MIMO communications. Typically, in SSK modulation, only one transmit antenna is activated during each symbol period. At the receiver, the index of the activated transmit antenna can be detected, when the channel state information is available. The SSK modulation can be implemented in conjunction with a conventional amplitude-phase modulation, forming the SM \\cite{LLY.PSM.11}. Similarly, the SM-MC proposed in this paper is the combination of a SSK modulation and a CSK modulation. Specifically, given a symbol transmitted, one of the transmitters releases a pulse of molecules, with the number of molecules determined also by the data symbol being transmitted. In detail, the transmit schematic diagram for our SM-MC system is depicted in Fig.~3. Hence, when the $j$-th space symbol is transmitted, with $j\\in\\{1,2,\\ldots, N\\}$, the transmit signal vector has a form of\n\\begin{align}\n\\mathbf{x}(t)=\\Big[0,0,\\ldots,\\underbrace{S_{m}}_{\\text{j-th transmitter}},\\ldots, 0,0\\Big]^{T},\n\\end{align}\nwhere only the $j$-th transmitter is activated to emit $S_{m}>0$ molecules, when the $m$-th concentration symbol $S_{m}$ is selected. The space symbol is denoted as $S^{j}$, when the $j$-th transmitter is activated. \nWe assume that the space symbols $S^{j}$ and the concentration symbols $S_{m}$ are independent of each other, solely depended on the input data stream. Then, we have\n\\begin{align}\n\\text{Pr}[S_{m}]=\\frac{1}{M},\\quad m\\in\\{0,1,\\ldots, M-1\\},\n\\end{align}\nand\n\\begin{align}\n\\text{Pr}[S^{j}]=\\frac{1}{N},\\quad j\\in\\{1,2,\\ldots, N\\}.\n\\end{align}\n\nFurthermore, if we express the SM symbols representing the combinations of the space and concentration symbols as $S^{j}_{m}$, we have\n\\begin{align}\n\\text{Pr}[S^{j}_{m}]=\\frac{1}{MN},\n\\begin{array}{c}\nj\\in\\{1,2,\\ldots, N\\},\\\\\nm\\in\\{0,1,\\ldots, M-1\\}.\\\\\n\\end{array}\n\\end{align}\n\nTherefore, the data rate of the SM-MC measured in bits per molecular symbol is given as\n\\begin{align}\nR_{\\sss \\text{SM}}=\\log_{2}N+\\log_{2}M,\n\\end{align}\nwhere both the values of $N$ and $M$ are assumed to be an integer power of $2$.\n\nBased on (17), the concentration signal observed at the $i$-th receiver in SM-MC systems is similar to (14), expressed as\n\\begin{align}\ny_{i,\\sss \\text{ SM}}(t)=S_{m}{h}_{ij}(t)+{I}_{i, \\sss \\text{ SM}}(t)+n_{i,\\sss \\text{ SM}}(t),\n\\end{align}\nwhere both $S_{m}{h}_{ij}(t)$ and $n_{i,\\sss \\text{ SM}}(t)$ are dependent on the current molecular symbol being received, representing the expected number of molecules received and the noise component, respectively, at time $t$, when the $j$-th transmitter is activated to emit a chemical impulse with $S_{m}$ molecules. However, unlike the case in MIMO-MC \\cite{Chae.MCMIMOproto.16}, the ${I}_{i, \\sss \\text{ SM}}(t)$ in (22) consists of only the ISI component that is resulted from a previous molecular symbol emitted by the $i$-th transmitter, since we assume that ILI only occurs with the current transmission, while ISI is only experienced from one previous emission by a paired transmitter. Let us denote the previous molecular symbol as $S^{\\bar{j}}_{\\bar{m}}$, then we have\n\\begin{align}\n{I}_{i, \\sss \\text{ SM}}(t)=\n\\begin{cases}\n\\ 0, & \\text{for $i\\neq j$}, \\\\\n\\ S_{\\bar{m}}{h}_{ii}(t+T_{s}),& \\text{for $i= \\bar{j}$.}\n\\end{cases}\n\\end{align}\n\nSimilar to (15), the distribution of the noise component in (22) can be expressed as\n\\begin{align}\nn_{i,\\sss \\text{ SM}}(t)\\sim \\mathcal{N} \\Big(\\mu_{ni,\\sss \\text{ SM}}(t),\\sigma^{2}_{ni,\\sss \\text{ SM}}(t)\\Big),\n\\end{align}\nassociated with\n\\begin{align}\n\\mu_{ni,\\sss \\text{ SM}}(t)=0,\\quad \\sigma^{2}_{ni,\\sss \\text{ SM}}(t)=\\frac{S_{m}{h}_{ij}(t)+{I}_{i, \\sss \\text{ SM}}(t)}{V_{\\sss \\text{RX}}}.\\nonumber\n\\end{align}\n\nFurthermore, based on (22), (23) and (24), the distribution of $y_{i,\\sss \\text{ SM}}(t)$ can be found, and expressed as\n\\begin{align}\ny_{i,\\sss \\text{ SM}}(t)\\hspace{-0.1cm}\\sim\\hspace{-0.05cm}\\mathcal{N}\\Big(\\hspace{-0.1cm}S_{m}{h}_{ij}(t)+{I}_{i, \\sss \\text{ SM}}(t),\\sigma^{2}_{ni,\\sss \\text{ SM}}(t)\\hspace{-0.05cm}\\Big).\n\\end{align}\n\n\\begin{table}[t]\n\t\\caption{Bits per pulse for different modulation schemes in MC} \n\t\\begin{tabular*}{9cm}{lclc} \n\t\t\\toprule \n\t\tScheme & Bits per pulse&Scheme & Bits per pulse \\\\ \n\t\t\\toprule \n\t\tBMoSK& 1&$N\\times N$ MIMO (OOK) & 2\\\\\n\t\tBCSK & 1&$2\\times2$ SM (BCSK)& 2 \\\\\n\t\tBSSK & 1&$8$SSK & 3\\\\\n\t\tOOK & 2&$2\\times2$ SM (QCSK)& 3\\\\\n\t\tQCSK & 2&$4\\times4$ SM (BCSK)& 3\\\\\n\t\tQSSK & 2&$16$SSK & 4\\\\ \n\t\t\\bottomrule \n\t\\end{tabular*} \n\\end{table} \nFrom the above analysis, we can readily realize that our proposed SM-MC transmission scheme is capable of mitigating the ILI experienced by the general MIMO-MC \\cite{Chae.MCMIMOproto.16}. Furthermore, the SM-MC scheme is more energy-efficient than the MIMO-MC. In MIMO-MC, when fixing a modulation scheme, increasing $N$ may transmit more bits per symbol but at the expense of higher energy consumption. By contrast, SM-MC and its special scheme of SSK-MC can obtain a logarithmic increase in bits per pulse transmission, which is contributed by the space modulation, with the increase of $N$, but without demanding any extra transmission energy. In order to highlight the energy efficiency achieved by the SM-MC and SSK-MC schemes, we compare the different modulation techniques of MC in Table~I. Note that in Table~I, we assume that the input data are independently uniformly distributed binary data. Hence, for example, when the OOK is used, transmitting $M$ bits requires about $M\/2$ pulses. Therefore, we have the average bits per pulse equaling 2. As another example, when the $N\\times N$ MIMO-MC using OOK is considered, it also has the average bits per pulse equaling 2.\n\n\\subsection{Signal Detection in SM-MC Systems}\nIn theory, the maximum-likelihood (ML) detection can be implemented by the SM-MC to jointly decode the index of the activated transmitter and the concentration symbol. When the memoryless receivers are taken into account, the ML detection based on (22) can be formulated as\n\\begin{align}\n<\\hat{j},\\hat{m}>=\\mathop{\\arg\\min}_{j\\in\\{1,2,\\dots,N\\},m\\in\\{0,1,\\dots,M-1\\}}\\left \\lVert\\mathbf{y}_{\\sss \\text{SM}}(t)-S_{m}\\mathbf{h}_{j}(t)\\right \\rVert^{2},\n\\end{align}\nwhere $\\hat{j}$ and $\\hat{m}$ are the estimated indices of the space and concentration symbols, respectively. Note that $S_{m}$ in SM-MC is always greater than zero, i.e., $S_{m}>0$, otherwise the space symbol is unable to be detected. However, (26) imposes a search complexity of $\\mathcal{O}(NM)$, when $NM$ is relatively large, hence it is not practical to deploy in MC systems, due to the constraint of the size and computing power of MC receivers. To this end, we suggest a low-complexity successive detection based on a scheme that has been commonly used in SM for wireless communications \\cite{JJ.optdetectSM.08,RZ.GPSM.13}. In this scheme, the CSK symbol is detected after the detection of the space symbol. To be more specific, we detect the index of the activated transmitter via the concentration comparison of the $N$ receivers based on the fact that the receiver paired with the activated transmitter is most likely to have the maximum concentration at the sampling time, as it is located with the minimum distance from its paired transmitter. Accordingly, the detection of the space symbol can be formulated as\n\\begin{align}\n\\hat{j}= \\mathop{\\arg\\max}_{j\\in\\{1,2,\\dots,N\\}}\\;{y_{j,\\sss \\text{ SM}}(t)}.\n\\end{align}\nAfter the detection of the space symbol, the $\\hat{j}$-th receiver's concentration can be used to detect the concentration symbol, which can be described as\n\\begin{align}\n\\hat{m}=\\mathop{\\arg\\min}_{m\\in\\{0,1,\\dots,M-1\\}}|{y}_{\\hat{j},\\sss \\text{ SM}}(t)-S_{m}h_{{\\hat{j}\\hat{j}}}(t)|^{2},\n\\end{align}\n\nFrom (27) and (28), we can readily know that the complexity of our proposed detector is $\\mathcal{O}(N+M)$. Provided that $N\\geqslant2$, $M\\geqslant2$, the complexity of the proposed detection scheme is lower than that of the ML detection of (26). It is worth noting that, once there is a transmitter activated, all the $N$ receivers can collect the molecules released. Hence, all the receivers can collaborate to make more reliable detection of the CSK symbol at the second stage. We refer to this detector as the EGC, described as\n\\begin{align}\n\\hat{m}_{\\sss \\text{EGC}}&=\\mathop{\\arg\\min}_{m\\in\\{0,1,\\dots,M-1\\}}\\left \\lVert\\mathbf{y}_{\\sss \\text{SM}}(t)-S_{m}\\mathbf{h}_{\\hat{j}}(t)\\right \\lVert^{2}.\n\\end{align}\nIt can be shown that the complexity of (29) combined with (27) is still $\\mathcal{O}(N+M)$, although the total number of computation of (29) is significantly higher than that of (28). However, the performance of the EGC-assisted detection can significantly outperform that of the detector of (28), which will be validated by the simulation results in Section IV.\n\\subsection{Error Performance Analysis of SM-MC}\nIn this subsection, we focus on the SER analysis of the successive detection, given by (27) and (29). Since the detection is dependent on the current molecular symbol $S^{j}_{m}$ and the previous one $S^{\\bar{j}}_{\\bar{m}}$, let us denote the correct detection probability of both space and concentration symbols as $P_{c,\\sss \\text{ SSK}}(S^{j}_{m}|S^{\\bar{j}}_{\\bar{m}})$ and $P_{c,\\sss \\text{ CSK}}(S^{j}_{m}|S^{\\bar{j}}_{\\bar{m}})$, respectively. Remembering that the $S^{j}_{m}$ and the $S^{\\bar{j}}_{\\bar{m}}$ are mutually independent, we can express the correct detection probability of space symbol as\n\\begin{align}\nP_{c,\\sss \\text{ SSK}}(S^{j}_{m}|S^{\\bar{j}}_{\\bar{m}})=\\prod_{i\\neq j}Q\\left(-\\frac{\\mu_{ji|S^{\\bar{j}}_{\\bar{m}}}}{\\sigma_{ji|S^{\\bar{j}}_{\\bar{m}}}}\\right),\n\\end{align}\nwhere $Q(\\cdot)$ is the Q-function. When given $i\\neq j$, $\\mu_{ji|S^{\\bar{j}}_{\\bar{m}}}$ and $\\sigma^{2}_{ji|S^{\\bar{j}}_{\\bar{m}}}$ are shown as\n\\begin{align}\n\\mu_{ji|S^{\\bar{j}}_{\\bar{m}}}\\hspace{-0.15cm}=\\hspace{-0.15cm}\n\\begin{cases}\n\\vspace{0.15cm} \n\\\n\\hspace{-0.225cm}S_{m}\\hspace{-0.075cm}\\big(h_{jj}\\hspace{-0.05cm}(t)\\hspace{-0.1cm}-\\hspace{-0.1cm}h_{ij}\\hspace{-0.05cm}(t)\\hspace{-0.06cm}\\big)\\hspace{-0.1cm}-\\hspace{-0.1cm}S_{\\bar{m}}h_{ii}\\hspace{-0.05cm}(t\\hspace{-0.1cm}+\\hspace{-0.1cm}T_{s}\\hspace{-0.05cm}),\\hspace{-0.3cm}&\\text{for $i= \\bar{j}$ and $j\\neq \\bar{j}$,}\\\\\n\\\n\\vspace{0.15cm} \\hspace{-0.225cm}S_{m}\\hspace{-0.075cm}\\big(h_{jj}\\hspace{-0.05cm}(t)\\hspace{-0.1cm}-\\hspace{-0.1cm}h_{ij}\\hspace{-0.05cm}(t)\\hspace{-0.06cm}\\big)\\hspace{-0.1cm}+\\hspace{-0.1cm}S_{\\bar{m}}h_{jj}\\hspace{-0.05cm}(t\\hspace{-0.1cm}+\\hspace{-0.1cm}T_{s}\\hspace{-0.05cm}),\\hspace{-0.3cm}&\\text{for $i\\neq \\bar{j}$ and $j= \\bar{j}$,}\\\\\n\\ \\hspace{-0.225cm}S_{m}\\hspace{-0.075cm}\\big(h_{jj}\\hspace{-0.05cm}(t)\\hspace{-0.1cm}-\\hspace{-0.1cm}h_{ij}\\hspace{-0.05cm}(t)\\hspace{-0.06cm}\\big), \\hspace{-0.3cm}&\\text{for $i\\neq \\bar{j}$ and $j\\neq \\bar{j}$.}\\nonumber\n\\end{cases}\n\\end{align}\n\\begin{align}\n\\sigma^{2}_{ji|S^{\\bar{j}}_{\\bar{m}}}\\hspace{-0.15cm}=\\hspace{-0.15cm}\n\\begin{cases}\n\\vspace{0.15cm}\n\\ \\hspace{-0.25cm}\\frac{S_{m}\\big(h_{jj}(t)+h_{ij}(t)\\big)+S_{\\bar{m}}h_{ii}(t+T_{s})}{V_{\\sss\\text{RX}}},&\\text{for $i= \\bar{j}$ and $j\\neq \\bar{j}$,}\\\\\n\\vspace{0.15cm}\n\\ \\hspace{-0.25cm}\\frac{S_{m}\\big(h_{jj}(t)+h_{ij}(t)\\big)+S_{\\bar{m}}h_{jj}(t+T_{s})}{V_{\\sss\\text{RX}}},&\\text{for $i\\neq \\bar{j}$ and $j= \\bar{j}$,}\\\\\n\\ \\hspace{-0.25cm}\\frac{S_{m}\\big(h_{jj}(t)+h_{ij}(t)\\big)}{V_{\\sss\\text{RX}}},&\\text{for $i\\neq \\bar{j}$ and $j\\neq \\bar{j}$.}\\nonumber\n\\end{cases}\n\\end{align}\n\n\\begin{IEEEproof}\n\tSee Appendix A.\n\\end{IEEEproof} \n\nSpecifically, when only the space symbol $S^{j}$ is transmitted, the SM-MC with $M=1$ degenerates to the SSK-MC and the detection in (27) is non-coherent. The detection performance of the SSK-MC is solely described by the expectation of $P_{c,\\sss \\text{ SSK}}(S^{j}|S^{\\bar{j}})$, where $S^{j}$ and $S^{\\bar{j}}$ are also mutually independent. Note that in SSK-MC, both $S^{j}$ and $S^{\\bar{j}}$ are integers. Correspondingly, considering that the transmitted symbols are uniform distribution, we have\n\\begin{align}\nP_{c,\\sss \\text{ SSK}}&=\\mathbb{E}\\Big[\\mathbb{E}\\big[P_{c,\\sss \\text{ SSK}}(S^{j}|S^{\\bar{j}})\\big]\\Big]\\nonumber\\\\\n&=\\sum_{j=1}^{N}\\sum_{\\bar{j}=1}^{N}P_{c,\\sss \\text{ SSK}}(S^{j}|S^{\\bar{j}})\\text{Pr}[S^{\\bar{j}}]\\text{Pr}[S^{j}]\\nonumber\\\\\n&=\\frac{1}{N^{2}}\\sum_{j=1}^{N}\\sum_{\\bar{j}=1}^{N}\\prod_{i\\neq j}Q\\left(\\frac{-\\mu_{ji|S^{\\bar{j}}}}{\\sigma_{ji|S^{\\bar{j}}}}\\right).\n\\end{align}\nNote that in (31), the subscript $m$ is omitted, since there is no concentration symbol transmitted in SSK-MC. Furthermore, the error probability can be easily derived as\n\\begin{align}\nP_{e,\\sss \\text{ SSK}}=1-P_{c,\\sss \\text{ SSK}}=1-\\frac{1}{N^{2}}\\sum_{j=1}^{N}\\sum_{\\bar{j}=1}^{N}\\prod_{i\\neq j}Q\\left(\\frac{-\\mu_{ji|S^{\\bar{j}}}}{\\sigma_{ji|S^{\\bar{j}}}}\\right).\n\\end{align}\n\nHaving considered the space symbol, the error rate of the CSK detection can be analyzed based on the estimate $\\hat{j}$. Here, we analyze the more general EGC-assisted detector of (29), which utilizes all receivers' observations to detect the CSK symbol. Given that a pulse of $S_{m}$ molecules is transmitted by the $j$-th transmitter and that the previously emitted molecular symbol is $S^{\\bar{j}}_{\\bar{m}}$, the probability of erroneous detection of the CSK symbol can be upper bounded by\n\\begin{align}\nP_{e,\\sss \\text{ CSK}}(S^{j}_{m}|S^{\\bar{j}}_{\\bar{m}})&\\leqslant \\sum_{m\\neq n}\\text{Pr}[S_{m}\\rightarrow S_{n}|S^{\\bar{j}}_{\\bar{m}}],\n\\end{align}\nwhere $\\text{Pr}[S_{m}\\rightarrow S_{n}|S^{\\bar{j}}_{\\bar{m}}]$ is also dependent on the detection of the space symbol. Hence, it can be expressed as\n\\begin{align}\n\\text{Pr}[S_{m}\\rightarrow S_{n}|S^{\\bar{j}}_{\\bar{m}}]=&\\sum_{\\hat{j}\\neq j}\\text{Pr}_{\\sss \\text{ SSK}}[\\hat{j}|S^{\\bar{j}}_{\\bar{m}}]\\ \\text{Pr}[S_{m}\\rightarrow S_{n}|S^{\\bar{j}}_{\\bar{m}},\\hat{j}]\\nonumber\\\\\n&+\\text{Pr}_{\\sss \\text{ SSK}}[j|S^{\\bar{j}}_{\\bar{m}}]\\ \\text{Pr}[S_{m}\\rightarrow S_{n}|S^{\\bar{j}}_{\\bar{m}},j],\n\\end{align}\nwhere $\\text{Pr}_{\\sss \\text{ SSK}}[\\hat{j}|S^{\\bar{j}}_{\\bar{m}}]$ and $\\text{Pr}_{\\sss \\text{ SSK}}[j|S^{\\bar{j}}_{\\bar{m}}]$ have been derived in (43) of Appendix~A for the detection of the SSK symbol. Hence, below we only focus on the unknown components, i.e., $\\text{Pr}[S_{m}\\rightarrow S_{n}|S^{\\bar{j}}_{\\bar{m}},\\hat{j}]$ and $\\text{Pr}[S_{m}\\rightarrow S_{n}|S^{\\bar{j}}_{\\bar{m}},j]$, in (34), which are given by (35) and (36), respectively, shown on the top of the next page. In these formulas, the $\\sigma^{2}_{ni,\\sss \\text{ SM}}(t)$ component has been defined in (24).\n\\begin{figure*}[t]\n\t\\begin{align}\n\t&\\text{Pr}[S_{m}\\rightarrow S_{n}|S^{\\bar{j}}_{\\bar{m}},\\hat{j}]=\\text{Pr}\\bigg[\\sum_{i=1}^{N}|{y}_{i,\\sss \\text{ SM}}(t)-S_{m}{h}_{i\\hat{j}}(t)|^{2}>\\sum_{i=1}^{N}|{y}_{i,\\sss \\text{ SM}}(t)-S_{n}{h}_{i\\hat{j}}(t)|^{2}\\Big|S^{\\bar{j}}_{\\bar{m}}\\bigg]\\nonumber\\\\\n\t&=\n\t\\begin{cases}\n\t\\ 1-Q\\left(\\frac{-S_{\\bar{m}}{h}_{\\bar{j}\\hat{j}}(t){h}_{\\bar{j}j}(t+T_{s})-S_{m}\\sum_{i=1}^{N}{h}_{i\\hat{j}}(t){h}_{ij}(t)+\\frac{(S_{m}+S_{n})}{2}\\sum_{i=1}^{N}{h}^{2}_{i\\hat{j}}(t)}{\\sqrt{\\sum_{i=1}^{N}{h}^{2}_{i\\hat{j}}(t)\\sigma^{2}_{ni,\\sss \\text{ SM}}(t)}}\\right),&\\text{for $S_{m}>S_{n}$}, \\\\\n\t\\ Q\\left(\\frac{-S_{\\bar{m}}{h}_{\\bar{j}\\hat{j}}(t){h}_{\\bar{j}j}(t+T_{s})-S_{m}\\sum_{i=1}^{N}{h}_{i\\hat{j}}(t){h}_{ij}(t)+\\frac{(S_{m}+S_{n})}{2}\\sum_{i=1}^{N}{h}^{2}_{i\\hat{j}}(t)}{\\sqrt{\\sum_{i=1}^{N}{h}^{2}_{i\\hat{j}}(t)\\sigma^{2}_{ni,\\sss \\text{ SM}}(t)}}\\right),&\\text{for $S_{m}\\sum_{i=1}^{N}|{y}_{i,\\sss \\text{ SM}}(t)-S_{n}{h}_{ij}(t)|^{2}\\Big|S^{\\bar{j}}_{\\bar{m}}\\bigg]\\nonumber\\\\\n\t&=\n\t\\begin{cases}\n\t\\ 1-Q\\left(\\frac{-S_{\\bar{m}}{h}_{\\bar{j}\\hat{j}}(t){h}_{\\bar{j}j}(t+T_{s})+\\frac{S_{n}-S_{m}}{2}\\sum_{i=1}^{N}{h}^{2}_{ij}(t)}{\\sqrt{\\sum_{i=1}^{N}{h}^{2}_{ij}(t)\\sigma^{2}_{ni,\\sss \\text{ SM}}(t)}}\\right), & \\text{for $S_{m}>S_{n}$}, \\\\\n\t\\ Q\\left(\\frac{-S_{\\bar{m}}{h}_{\\bar{j}\\hat{j}}(t){h}_{\\bar{j}j}(t+T_{s})+\\frac{S_{n}-S_{m}}{2}\\sum_{i=1}^{N}{h}^{2}_{ij}(t)}{\\sqrt{\\sum_{i=1}^{N}{h}^{2}_{ij}(t)\\sigma^{2}_{ni,\\sss \\text{ SM}}(t)}}\\right), & \\text{for $S_{m}0$ should be satisfied. Eq. (40) implies that SNR is merely dependent on the molecule number of molecules released in average by a chemical impulse, the transceiver distance $d$ and the volume of receiver $V_{\\sss {\\text{RX}}}$. Similarly, the SNR of each link in the $N\\times N$ MIMO-MC with $M$-ary CSK can be defined as\n\\begin{align}\n{\\text{SNR}}=\\bigg(\\frac{3}{2\\pi e}\\bigg)^{\\frac{3}{2}}\\frac{V_{\\sss {\\text{RX}}}}{NMd^{3}}\\sum_{m=0}^{M-1}S_{m}.\n\\end{align}\nNote that (41) is $1\/N$ of the (40), indicating that all the transmitters are activated to emitted chemical pulses in each symbol duration, when $S_{m}>0$. Specifically, when the OOK modulation is used in the $N\\times N$ MIMO-MC, only $N\/2$ transmitters are activated in average to emitted chemical pulses with $S_{1}$ molecules during each symbol period. Hence, the SNR of each link in the $N\\times N$ MIMO-MC is then defined as\n\\begin{align}\n{\\text{SNR}}=\\bigg(\\frac{3}{2\\pi e}\\bigg)^{\\frac{3}{2}}\\frac{2V_{\\sss {\\text{RX}}}}{Nd^{3}}S_{1},\n\\end{align}\nwhere $S_{1}$ is the number of molecules for transmitting bit ``1''.\n\n=When the transceiver distance $d$ is fixed, the same SNR implies the same molecular energy consumption. We can make relatively fair SER comparison between different modulation schemes of MC with the same transmission rate, under the same energy consumption. In order to achieve this, in our simulation for SM-MC with BCSK, we set $S_{1}=2S_{0}$. By contrast, for SISO-MC with QCSK, we have $S_{0}=0$ and $S_{3}=\\frac{3}{2}S_{2}=3S_{1}$. The other system parameters are presented in Table~II. Note that in this section, the SSK-MC with the $2\\times2$ and $4\\times4$ MIMO architectures, which for brevity are referred to as the BSSK-MC and QSSK-MC.\n\\begin{figure}[tb]\n\t\\centering\n\t\\includegraphics[width=3.5in]{Fig_perform1SSK.eps}\n\t\\caption{SER comparison between BSSK-MC and QSSK-MC, when receiver separation distance $r=12.5\\ \\rm{\\mu m}$ and different symbol durations $T_{s}=0.1\\ \\rm{s},0.2\\ \\rm{s},0.8\\ \\rm{s}$ are assumed.}\n\t\\label{fig_4}\n\\end{figure}\n\nFig. 4 shows the analytical and simulation results for the SER of the BSSK-MC and QSSK-MC, when different symbol durations are considered. Clearly, the analytical SER agrees well with the corresponding simulated SER. In the low SNR region, the SER of QSSK-MC is much higher than that of the BSSK-MC, while in the high SNR region, the SER performance of QSSK-MC is becomes better than that of the BSSK-MC, when $T_{s}=0.1\\ \\rm{s}$ and $T_{s}=0.2\\ \\rm{s}$. By contrast, when $T_{s}=0.8\\ \\rm{s}$, the SER of BSSK-MC is always lower than that of QSSK-MC in the whole SNR region considered. However, the gap between their SER performance decreases with the increase of SNR. As shown in Fig. 4, the SER performance of BSSK-MC and QSSK-MC is highly dependent on the symbol duration $T_{s}$ . When $T_{s}$ increases, the SER reduces due to the fact that the ISI reduces, as $T_{s}$ increases.\n\\begin{figure}[tb]\n\t\\centering\n\t\\includegraphics[width=3.5in]{Fig_perform3SSK.eps}\n\t\\caption{SER comparison between BSSK-MC and QSSK-MC, when $T_{s}=0.5\\ \\rm{s}$, and different receiver separation distances are considered.}\n\t\\label{fig_5}\n\\end{figure}\n\nFig. 5 manifests the analytical and simulation SER of BSSK-MC and QSSK-MC, when different transceiver separation distances are considered. Again, the simulation and analytical results match well, which hence validate our theoretical analysis. As shown in the figure, for a given transceiver separation distance, there is a gap between the SER performance of BSSK-MC and that of QSSK-MC in the low SNR region, with BSSK-MC always outperforming the QSSK-MC. However, as the SNR increases, the SER curves of BSSK-MC and QSSK-MC converge and become nearly the same at high SNR. Therefore, when sufficient source of information molecules, i.e., SNR, is available, higher throughput can be attained by utilizing higher order modulation schemes in the SSK-MC, while achieving the required error performance.\n\\begin{figure}[tb]\n\t\\centering\n\t\\includegraphics[width=3.5in]{Fig_perform2SM.eps}\n\t\\caption{SER comparison between the $2\\times2$ SM-MC and $4\\times4$ SM-MC, when receiver separation distance $r=10\\ \\rm{\\mu m}$ and different symbol durations $T_{s}=0.15\\ \\rm{s},0.3\\ \\rm{s},1\\ \\rm{s}$ are assumed.}\n\t\\label{fig_6}\n\\end{figure}\n\n\nFig. 6 depicts the theoretical SER upper bound and the simulated SER of the SM-MC, when the $2\\times2$ and the $4\\times4$ MIMO architectures are respectively considered. In this figure, the transceiver separation distance is fixed to $r=10\\ \\rm{\\mu m}$, whilst the symbol durations $T_{s}$ is set to $0.15\\ \\rm{s}$, $0.3\\ \\rm{s}$ or $1\\ \\rm{s}$. The results show that the upper bound is tight, which becomes tighter as the SNR increases. As shown in Fig. 6, in the low SNR regime, the SER of the $4\\times4$ SM-MC is much higher than that of the $2\\times2$ SM-MC. By contrast, in the high SNR regime, the SER of the $4\\times4$ SM-MC is better than that of the $2\\times2$ SM-MC, when $T_{s}=0.15\\ \\rm{s}$ or $T_{s}=0.3\\ \\rm{s}$. Furthermore, when $T_{s}=1\\ \\rm{s}$, the SER of the BSSK-MC is always lower than that of the QSSK-MC in the SNR regime considered. The explanation for these results are similar to that for the results shown in Fig. 4.\n\\begin{figure}[tb]\n\t\\centering\n\t\\includegraphics[width=3.5in]{Fig_perform4SM.eps}\n\t\\caption{SER comparison between the $2\\times2$ SM-MC and $4\\times4$ SM-MC, when $T_{s}=1\\ \\rm{s}$, and $r=8\\ \\rm{\\mu m},10\\ \\rm{\\mu m},12\\ \\rm{\\mu m}$ are assumed.}\n\t\\label{fig_7}\n\\end{figure}\n\nFig. 7 also demonstrates the theoretical SER upper bound and the simulation SER results of the SM-MC, where the symbol duration is set to $T_{s}=1\\ \\rm{s}$, while various transceiver separation distances are considered. Similar to Fig. 6, the SER upper bounds match well with their corresponding results obtained from simulations. There is also an SER gap between the $2\\times2$ SM-MC and its $4\\times4$ counterpart. However, when the transceiver separation distance increases, the SER difference evaluated in $\\rm{dB}$ reduces. We should note that when the link distance $d$ and the symbol duration $T_{s}$ are fixed, the SER performance of SM-MC is mainly determined by the separation distance $r$ between the receivers.\n\\begin{figure}[tb]\n\t\\centering\n\t\\includegraphics[width=3.5in]{Fig_perform5SM.eps}\n\t\\caption{SER comparison of the SM-MC and SSK-MC with that of the MIMO-MC using OOK modulation and SISO-MC using QCSK modulation, when $T_{s}=0.2\\ \\rm{s}$, and $r=10\\ \\rm{\\mu m}$ or $15\\ \\rm{\\mu m}$.}\n\t\\label{fig_8}\n\\end{figure}\n\\begin{figure}[tb]\n\t\\centering\n\t\\includegraphics[width=3.5in]{Fig_perform6SM.eps}\n\t\\caption{SER performance of the SM-MC with EGC-assisted and SC-assisted detection schemes, when $T_{s}=1\\ \\rm{s}$, and $r=10\\ \\rm{\\mu m},12.5\\ \\rm{\\mu m}$ or $15\\ \\rm{\\mu m}$.}\n\t\\label{fig_9}\n\\end{figure}\n\nTo verify the effectiveness of our proposed SSK-MC and SM-MC, we compare the simulated SER performance of the SISO-MC with QCSK, $2 \\times 2$ MIMO-MC with OOK, $4 \\times 4$ SSK-MC and the $2 \\times 2$ SM-MC in Fig. 8. All of them have the same transmission rate of 2 bits per symbol. As shown in Fig. 8, when $r=15\\ \\rm{\\mu m}$, the QSSK-MC outperforms all the other schemes and achieves the best SER performance. Both Both the SSK-MC and SM-MC benefit from the employment of space modulation, making them significantly surpass the $2 \\times 2$ MIMO-MC in the high SNR regime. The SISO-MC with QCSK attains the worst SER performance among the schemes considered. However, the situation is completely different when the separation distance is reduced to $r=10\\ \\rm{\\mu m}$. In this case, the SISO-MC with QCSK achieves the lowest SER, and outperforms all the other MIMO schemes. As shown in Fig. 8, the $2 \\times 2$ MIMO-MC with OOK has similar SER performance as the QSSK-MC and the $2 \\times 2$ SM-MC when the SNR is between $0\\ \\rm{dB}$ and $12\\ \\rm{dB}$. When for the increasing SNR, its SER performance becomes poorer than that of the other schemes.\n\nThe results in Fig. 8 implies that the separation distance $r$ is important for the achievable SER performance, especially when comparing the modulation schemes in SISO-MC. When $r$ is appropriately selected, the MIMO-MC schemes are capable of outperforming SISO-MC with QCSK, owning to the significant ILI reduction in the MIMO-MC schemes, which is beneficial to transmitting information in the space domain. By contrast, if $r$ is not sufficiently large, the advantage of using space symbol vanishes due to the high ILI.\n\nIn Fig. 9, we compare the SER performance of the SC-assisted SM-MC of (28) with that of the EGC-assisted SM-MC of (29). From Fig. 9, we observe that the SER performance of EGC-assisted SM-MC is better than that of SC-assisted SM-MC, when there is only a single receiver utilized to decode the concentration symbol. For the $4 \\times 4$ SM-MC with $r=10\\ \\rm{\\mu m}$, EGC-assisted SM-MC slightly outperforms SC-assisted SM-MC, when the SNR is lower than about $16 \\ \\rm{dB}$, while their performance converges at high SNR. When $r=12.5\\ \\rm{\\mu m}$, and the $2 \\times 2$ and $4 \\times 4$ SM-MC are considered, we observe that the EGC-assisted SM-MC always outperforms its SC-assisted counterpart. Additionally, when the same combining scheme is employed at the receivers, a small-scale SM-MC with fewer links outperforms a large-scale SM-MC with more links. Fig.9 also shows, when $2 \\times 2$ SM-MC with $r=15\\ \\rm{\\mu m}$ is considered, the SER performance of EGC-assisted SM-MC becomes increasingly better than that of SC-assisted SM-MC, and approximately $1.5\\ \\rm{dB}$ gain can be obtained at the SER of $10^{-3}$.\n\nBased on the aforementioned results, in order to improve the SER performance of the SSK-MC and that of SM-MC, the following approaches may be applied. Firstly, the symbol duration $T_{s}$ may be increased to reduce ISI, which hence improves the SER of the SSK-MC and SM-MC systems, but at the cost of the transmission rate. Secondly, the separation distance between adjacent transmitters and that between adjacent receivers may be increased for mitigating the ILI. This is an effective method to improve the SER performance provided that tehre are spaces for transmitter and receiver deployment. Furthermore, the number of molecules per pulse maybe increased, which increases SNR, and therefore enhances the reliability of the information transmission in SSK-MC and the SM-MC system, but at the cost of energy consumption.\n\n\\section{Conclusions and Future Work}\nIn this paper, we have proposed a SM-MC scheme and its special case of SSK-MC by introducing new degrees of freedom from the spatial domain to transmit more information in MC. Our studies showed that when the separation distance between transmit\/receive nanomachines is appropriate, additional information conveyed in the space domain is achievable, and our proposed SM-MC and SSK-MC schemes are capable of outperforming the conventional SISO-MC and MIMO-MC scheme. The reason behind is that our proposed SM-MC can effectively eliminate the ILI, which usually severely affects the performance of the conventional MIMO-MC. Moreover, it is an energy-efficient system where only one nanomachine transmitter is activated during each symbol period, befitting the twofold goal in small scale communication systems of low energy consumption and low-complexity.\n\nIn this paper, we have also presented theoretical analysis and some simulation results under the assumption of the homogeneous propagation medium as done in the literature. This may not accord with a realistic situation where various medium exist in the communication channel. In our future work, we will consider the diffusion based MC with more realistic channel modelling, as well as the design of corresponding MIMO-MC transceivers.\n\n\\begin{appendices}\n\\section{ }\nLet us denote $\\text{Pr}_{\\sss \\text{ SSK}}[i|S^{\\bar{j}}_{\\bar{m}}]$ as the probability that the $i$-th receiver senses the maximum molecular concentration, which can be expressed as\n\t\\begin{align}\n\t\\text{Pr}_{\\sss \\text{ SSK}}[i|S^{\\bar{j}}_{\\bar{m}}]=\\text{Pr}\\Big[\\hspace{-0.05cm}\\big(y_{i,\\sss \\text{ SM}}(t)\\hspace{-0.1cm}>\\hspace{-0.05cm}y_{1,\\sss \\text{ SM}}(t)\\hspace{-0.05cm}\\big)\\hspace{-0.1cm}\\cap\\hspace{-0.1cm}\\cdots\\hspace{-0.1cm}\\cap\\hspace{-0.1cm}\\big(y_{i,\\sss \\text{ SM}}(t)\\hspace{-0.1cm}>\\hspace{-0.05cm}y_{N,\\sss \\text{ SM}}(t)\\hspace{-0.05cm}\\big)\\hspace{-0.05cm}\\Big],\n\t\\end{align}\nwhen given that the previously sent molecular symbol is $S^{\\bar{j}}_{\\bar{m}}$. In order to derive the error detection probability of the space symbol, let us first derive the correct detection probability of the space symbol, when $S_{m}$ molecules are released from the $j$-th transmitter. Based on (43), this probability can be expressed as\n\\begin{align}\nP_{c,\\sss \\text{ SSK}}(S^{j}_{m}|S^{\\bar{j}}_{\\bar{m}})&=\\text{Pr}_{\\sss \\text{ SSK}}[j|S^{\\bar{j}}_{\\bar{m}}]\\nonumber\\\\\n&=\\text{Pr}\\hspace{-0.05cm}\\Big[\\hspace{-0.075cm}\\big(y_{j,\\sss \\text{ SM}}(t)\\hspace{-0.1cm}>\\hspace{-0.1cm}y_{1,\\sss \\text{ SM}}(t)\\hspace{-0.075cm}\\big)\\hspace{-0.1cm}\\cap\\hspace{-0.1cm}\\cdots\\hspace{-0.1cm}\\cap\\hspace{-0.1cm}\\big(y_{j,\\sss \\text{ SM}}(t)\\hspace{-0.1cm}>\\hspace{-0.1cm}y_{N,\\sss \\text{ SM}}(t)\\hspace{-0.075cm}\\big)\\hspace{-0.075cm}\\Big],\n\\end{align}\nwhich is the probability that the $j$-th receiver senses the maximum concentration. When the reception processes of the receiver nanomachines are assumed to be independent, (44) can be expressed as\n\\begin{align}\nP_{c,\\sss \\text{ SSK}}(S^{j}_{m}|S^{\\bar{j}}_{\\bar{m}})=\\prod_{i\\neq j}\\text{Pr}[y_{j,\\sss \\text{ SM}}(t)-y_{i,\\sss \\text{ SM}}(t)>0|S^{\\bar{j}}_{\\bar{m}}].\n\\end{align}\n\nFurthermore, when given $S^{\\bar{j}}_{\\bar{m}}$, $y_{j,\\sss \\text{ SM}}(t)-y_{i,\\sss \\text{ SM}}(t)$ can be written as\n\\begin{align}\ny_{ji|S^{\\bar{j}}_{\\bar{m}}}=&S_{m}\\big(h_{jj}(t)-h_{ij}(t)\\big)+I_{j,\\sss \\text{ SM}}(t)-I_{i,\\sss \\text{ SM}}(t)\\nonumber\\\\&+n_{j,\\sss \\text{ SM}}(t)-n_{i,\\sss \\text{ SM}}(t).\n\\end{align}\nAccording to (25), $y_{ji|S^{\\bar{j}}_{\\bar{m}}}$ follows the Normal distribution of\n\\begin{align}\ny_{ji|S^{\\bar{j}}_{\\bar{m}}}\\sim \\mathcal{N} \\left(\\mu_{ji|S^{\\bar{j}}_{\\bar{m}}},\\sigma^{2}_{ji|S^{\\bar{j}}_{\\bar{m}}}\\right).\n\\end{align}\nWhen given $i\\neq j$, and according to \\cite{Pat.NormalDist.96}, $\\mu_{ji|S^{\\bar{j}}_{\\bar{m}}}$ and $\\sigma^{2}_{ji|S^{\\bar{j}}_{\\bar{m}}}\\hspace{-0.15cm}$ are given by\n\\begin{align}\n\\mu_{ji|S^{\\bar{j}}_{\\bar{m}}}\\hspace{-0.15cm}=\\hspace{-0.15cm}\n\\begin{cases}\n\\vspace{0.15cm} \n\\\n\\hspace{-0.225cm}S_{m}\\hspace{-0.075cm}\\big(h_{jj}\\hspace{-0.05cm}(t)\\hspace{-0.1cm}-\\hspace{-0.1cm}h_{ij}\\hspace{-0.05cm}(t)\\hspace{-0.06cm}\\big)\\hspace{-0.1cm}-\\hspace{-0.1cm}S_{\\bar{m}}h_{ii}\\hspace{-0.05cm}(t\\hspace{-0.1cm}+\\hspace{-0.1cm}T_{s}\\hspace{-0.05cm}),\\hspace{-0.3cm}&\\text{for $i= \\bar{j}$ and $j\\neq \\bar{j}$,}\\\\\n\\\n\\vspace{0.15cm} \\hspace{-0.225cm}S_{m}\\hspace{-0.075cm}\\big(h_{jj}\\hspace{-0.05cm}(t)\\hspace{-0.1cm}-\\hspace{-0.1cm}h_{ij}\\hspace{-0.05cm}(t)\\hspace{-0.06cm}\\big)\\hspace{-0.1cm}+\\hspace{-0.1cm}S_{\\bar{m}}h_{jj}\\hspace{-0.05cm}(t\\hspace{-0.1cm}+\\hspace{-0.1cm}T_{s}\\hspace{-0.05cm}),\\hspace{-0.3cm}&\\text{for $i\\neq \\bar{j}$ and $j= \\bar{j}$,}\\\\\n\\ \\hspace{-0.225cm}S_{m}\\hspace{-0.075cm}\\big(h_{jj}\\hspace{-0.05cm}(t)\\hspace{-0.1cm}-\\hspace{-0.1cm}h_{ij}\\hspace{-0.05cm}(t)\\hspace{-0.06cm}\\big), \\hspace{-0.3cm}&\\text{for $i\\neq \\bar{j}$ and $j\\neq \\bar{j}$,}\n\\end{cases}\n\\end{align}\n\\begin{align}\n\\sigma^{2}_{ji|S^{\\bar{j}}_{\\bar{m}}}\\hspace{-0.15cm}=\\hspace{-0.15cm}\n\\begin{cases}\n\\vspace{0.15cm}\n\\ \\hspace{-0.25cm}\\frac{S_{m}\\big(h_{jj}(t)+h_{ij}(t)\\big)+S_{\\bar{m}}h_{ii}(t+T_{s})}{V_{\\sss\\text{RX}}},&\\text{for $i= \\bar{j}$ and $j\\neq \\bar{j}$,}\\\\\n\\vspace{0.15cm}\n\\ \\hspace{-0.25cm}\\frac{S_{m}\\big(h_{jj}(t)+h_{ij}(t)\\big)+S_{\\bar{m}}h_{jj}(t+T_{s})}{V_{\\sss\\text{RX}}},&\\text{for $i\\neq \\bar{j}$ and $j= \\bar{j}$,}\\\\\n\\ \\hspace{-0.25cm}\\frac{S_{m}\\big(h_{jj}(t)+h_{ij}(t)\\big)}{V_{\\sss\\text{RX}}},&\\text{for $i\\neq \\bar{j}$ and $j\\neq \\bar{j}$,}\n\\end{cases}\n\\end{align}\nrespectively. Therefore, we have\n\\begin{align}\n\\text{Pr}[y_{ji|S^{\\bar{j}}_{\\bar{m}}}\\hspace{-0.1cm}>\\hspace{-0.1cm}0]\\hspace{-0.1cm}&=\\hspace{-0.2cm}\\int_{0}^{+\\infty}\\hspace{-0.5cm}\\frac{1}{\\sqrt{2\\pi}\\sigma_{ji|S^{\\bar{j}}_{\\bar{m}}}}\\exp\\hspace{-0.1cm}\\left(\\hspace{-0.15cm}-\\frac{(y_{ji|S^{\\bar{j}}_{\\bar{m}}}-\\mu_{ji|S^{\\bar{j}}_{\\bar{m}}})^{2}}{2\\sigma^{2}_{ji|S^{\\bar{j}}_{\\bar{m}}}}\\hspace{-0.125cm}\\right)\\hspace{-0.125cm}dy_{ji|S^{\\bar{j}}_{\\bar{m}}}\\nonumber\\\\\n&=Q\\left(-\\frac{\\mu_{ji|S^{\\bar{j}}_{\\bar{m}}}}{\\sigma_{ji|S^{\\bar{j}}_{\\bar{m}}}}\\right).\n\\end{align}\nFinally, when substituting (50) into (45), the probability of correct detection of the space symbol $j$ can be formulated as\n\\begin{align}\nP_{c,\\sss \\text{ SSK}}(S^{j}_{m}|S^{\\bar{j}}_{\\bar{m}})=\\prod_{i\\neq j}Q\\left(-\\frac{\\mu_{ji|S^{\\bar{j}}_{\\bar{m}}}}{\\sigma_{ji|S^{\\bar{j}}_{\\bar{m}}}}\\right).\n\\end{align}\n\n\n\\end{appendices}\n\\bibliographystyle{IEEEtran} \n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\n The classification problem of finite dimensional algebras is one of the important problems of modern algebra. So far two approaches are known to the solution of the problem. One of them is structural (basis free, invariant) approach. For instance, the classification of finite dimensional simple and semi-simple associative algebras by Wedderburn and simple and semi-simple Lie algebras by Cartan are examples of such approach. But it is observed that this approach becomes more difficult when one considers more general types of algebras. Another approach to the solution of the problem is coordinate based approach (see \\cite{B2015,D2003,E2004,G2011,P2011} for the latest results). These two approaches somehow are complementary to each other.\n \n In two-dimensional case a complete classification, by basis free approach, is stated in \\cite{P2000} over any basic field. In this paper we follow the coordinate (basis, structure constants) based approach, we consider such problem for two-dimensional algebras over any algebraically closed field $\\mathbf{F}$. We provide a list of algebras, given by their matrices of structure constants (MSC), such that any 2-dimensional algebra is isomorphic to only one algebra from the list. A similar result has been stated in \\cite{G2011}. The main difference between these two results is as follows. In \\cite{G2011} the authors can state the existence only whereas the uniqueness can not be guaranteed. Moreover, the approach to the problem followed in the present paper is totally different that in \\cite{G2011}. Some details of our approach for arbitrary dimensional case have been given in \\cite{B2015}.\n \n The paper is organized as follows. First we study the problem over algebraically closed fields of characteristics not 2 and 3, then the solution is given for algebras over algebraically closed fields of characteristics 2 and the last subsection of paper contains the result over algebraically closed fields of characteristics 3. In each of these cases we present the list of algebras via their matrices of structure constants.\n \n\\section{Classification of two-dimensional algebras}\n\nTo classify the main part of two-dimensional algebras we use a particular case of the following result from \\cite{B2015}. Let $n$, $m$ be any natural numbers, $\\tau: (G,V)\\rightarrow V$ be a fixed linear algebraic representation of an algebraic subgroup $G$ of $GL(m,\\mathbf{F})$ on the $n$-dimensional vector space $V$ over $\\mathbf{F}$. Assume that there exists a nonempty $G$-invariant subset $V_0$ of $V$ and an algebraic map $P: V_0\\rightarrow G$ such that\n\\begin{center}\n\t$P(\\tau( g,\\mathbf{v}))=P(\\mathbf{v}) g^{-1}$,\n\\end{center}\nwhenever $\\mathbf{v}\\in V_0$ and $g\\in G$. The following result holds true \\cite{B2015}.\n\n{\\bf Theorem.} \\textit{ Elements $\\mathbf{u},\\mathbf{v}\\in V_0$ are $G$-equivalent, that is $\\mathbf{u}=\\tau( g,\\mathbf{v})$ for some $g\\in G$, if and only if \\[\\tau(\n\tP(\\mathbf{u}),\\mathbf{u})=\\tau( P(\\mathbf{v}),\\mathbf{v}).\\]}\n\nLet $\\mathbf{A}$ be any 2 dimensional algebra over $\\mathbf{F}$ with multiplication $\\cdot$ given by a bilinear map $(\\mathbf{u},\\mathbf{v})\\mapsto \\mathbf{u}\\cdot \\mathbf{v}$ whenever $\\mathbf{u}, \\mathbf{v}\\in \\mathbf{A}$. If $e=(e^1,e^2)$ is \nbasis for $\\mathbf{A}$ as a vector space over $\\mathbf{F}$ then one can represent this bilinear map by a matrix \\[A=\\left(\n\\begin{array}{cccc}\nA^{1}_{1,1} & A^{1}_{1,2}& A^{1}_{2,1}& A^{1}_{2,2} \\\\\nA^{2}_{1,1} & A^{2}_{1,2}& A^{2}_{2,1}& A^{2}_{2,2} \\\\\n\\end{array}\n\\right) \\in Mat(2\\times 4;\\mathbf{F})\\] such that \\[\\mathbf{u}\\cdot \\mathbf{v}=eA(u\\otimes v)\\] for any $\\mathbf{u}=eu,\\mathbf{v}=ev,$\nwhere $u=(u_1, u_2),$ and $v=(v_1, v_2)$ are column coordinate vectors of $\\mathbf{u}$ and $\\mathbf{v},$ respectively, $(u\\otimes v)=(u_1v_1,u_1v_2,u_2v_1,u_2v_2)$, $e^i\\cdot e^j=A^{1}_{i,j}e^1+A^{2}_{i,j}e^2$ whenever $i,j=1,2$.\nSo the algebra $\\mathbf{A}$ is presented by the matrix $A\\in Mat(2\\times 4;\\mathbf{F})$ (called the matrix of MSC of $\\mathbf{A}$ with respect to the basis $e$).\n\nIf $e'=(e'^1,e'^2)$ is also another basis for $\\mathbf{A}$, $g\\in G=GL(2,\\mathbf{F})$, $e'g=e$ and $\\mathbf{u}\\cdot \\mathbf{v}=e'B(u'\\otimes v')$, where $\\mathbf{u}=e'u',\\mathbf{v}=e'v'$, then\n\\[\\mathbf{u}\\cdot \\mathbf{v}=eA(u\\otimes v)=e'B(u'\\otimes v')=eg^{-1}B(gu\\otimes gv)=eg^{-1}B(g\\otimes g)(u\\otimes v)\\] as far as $\\mathbf{u}=eu=e'u'=eg^{-1}u',\\mathbf{v}=ev=e'v'=eg^{-1}v'$.\nTherefore the equality\n\\[B=gA(g^{-1})^{\\otimes 2}\\] is valid, where for $g^{-1}=\\left(\\begin{array}{cccc} \\xi_1& \\eta_1\\\\ \\xi_2& \\eta_2\\end{array}\\right)$ one has\n\\[(g^{-1})^{\\otimes 2}=g^{-1}\\otimes g^{-1}=\\left(\\begin{array}{cccc} \\xi _1^2 & \\xi _1 \\eta _1 & \\xi _1 \\eta _1 & \\eta _1^2 \\\\\n\\xi _1 \\xi _2 & \\xi _1 \\eta _2 & \\xi _2 \\eta _1 & \\eta _1 \\eta _2 \\\\\n\\xi _1 \\xi _2 & \\xi _2 \\eta _1 & \\xi _1 \\eta _2 & \\eta _1 \\eta _2 \\\\\n\\xi _2^2 & \\xi _2 \\eta _2 & \\xi _2 \\eta _2 & \\eta _2^2\n\\end{array}\n\\right).\\]\n\n{\\bf Definition.} Two-dimensional algebras $\\mathbf{A}$, $\\mathbf{B}$, given by\n\ttheir matrices of structural constants $A$, $B$, are said to be isomorphic if $B=gA(g^{-1})^{\\otimes 2}$ holds true for some $g\\in GL(2,\\mathbf{F})$.\n\nNote that the following identities \\begin{equation} \\label{2.1} Tr_1(gA(g^{-1})^{\\otimes 2})=Tr_1(A)g^{-1},\\ \\ Tr_2(gA(g^{-1})^{\\otimes 2})=Tr_2(A)g^{-1}\\end{equation} hold true whenever $A\\in Mat(2\\times 4,\\mathbf{F}), g\\in GL(2,\\mathbf{F})$, where \\[Tr_1(A)=(A^1_{1,1}+A^2_{2,1}, A^1_{1,2}+A^2_{2,2}),\\ Tr_2(A)=(A^1_{1,1}+A^2_{1,2}, A^1_{2,1}+A^2_{2,2})\\] are row vectors.\n\nWe divide $Mat(2\\times 4,\\mathbf{F})$ into the following five disjoint subsets:\n\\begin{itemize}\n\t\\item[1.] All $A$ for which the system $\\{Tr_1(A),Tr_2(A)\\}$ is linear independent.\n\t\\item[2.] All $A$ for which the system $\\{Tr_1(A),Tr_2(A)\\}$ is linear dependent and $Tr_1(A),Tr_2(A)$ are nonzero vectors.\n\t\\item[3.] All $A$ for which $Tr_1(A)$ is nonzero vector and $Tr_2(A)=(0,0)$.\n\t\\item[4.] All $A$ for which $Tr_1(A)=(0,0)$ and $Tr_2(A)$ is nonzero vector.\n\t\\item[5.] All $A$ for which$Tr_1(A)=Tr_2(A)=(0,0)$.\n\\end{itemize}\n\nDue to (1) it is clear that algebras with matrices from these different classes can not be isomorphic. We deal with each of these subsets separately.\nFurther, for the simplicity, we use the notation \\[A=\\left(\\begin{array}{cccc} \\alpha_1 & \\alpha_2 & \\alpha_3 &\\alpha_4\\\\ \\beta_1 & \\beta_2 & \\beta_3 &\\beta_4\\end{array}\\right),\\] where\n$\\alpha_1, \\alpha_2, \\alpha_3, \\alpha_4, \\beta_1, \\beta_2, \\beta_3, \\beta_4$ stand for any elements of $\\mathbf{F}$.\n\nAs it has been stated above we consider three cases\ncorresponding to $Char(\\mathbf{F})\\neq 2,3,$ $Char(\\mathbf{F})=2$ and $Char(\\mathbf{F})=3,$ respectively.\n\n\\subsection{Characteristic not $2$, $3$ case}\n\n{\\bf Theorem 1.} \\textit{Over an algebraically closed field $\\mathbf{F}$ with the characteristic not 2 and 3, any non-trivial 2-dimensional algebra is isomorphic to only one of the following listed, by their matrices of structure constants, algebras:\n\t\\[A_{1}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1 & \\alpha_2 &\\alpha_2+1 & \\alpha_4 \\\\\n\t\\beta_1 & -\\alpha_1 & -\\alpha_1+1 & -\\alpha_2\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\alpha_2, \\alpha_4, \\beta_1)\\in \\mathbf{F}^4,\\]\n\t\\[A_{2}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1 & 0 & 0 & 1 \\\\\n\t\\beta _1& \\beta _2& 1-\\alpha_1&0\n\t\\end{array}\\right)\\simeq \\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1 & 0 & 0 & 1 \\\\\n\t-\\beta _1& \\beta _2& 1-\\alpha_1&0\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1, \\beta_2)\\in \\mathbf{F}^3,\\]\n\t\\[A_{3}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t0 & 1 & 1 & 0 \\\\\n\t\\beta _1& \\beta _2 & 1&-1\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\beta_1, \\beta_2)\\in \\mathbf{F}^2,\\]\n\t\\[A_{4}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t\\alpha _1 & 0 & 0 & 0 \\\\\n\t0 & \\beta _2& 1-\\alpha _1&0\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_2)\\in \\mathbf{F}^2,\\]\n\t\\[A_{5}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1& 0 & 0 & 0 \\\\\n\t1 & 2\\alpha_1-1 & 1-\\alpha_1&0\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F},\\]\n\t\\[A_{6}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1 & 0 & 0 & 1 \\\\\n\t\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\t\\end{array}\\right)\\simeq \\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1 & 0 & 0 & 1 \\\\\n\t-\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1)\\in \\mathbf{F}^2,\\]\n\t\\[A_{7}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t0 & 1 & 1 & 0 \\\\\n\t\\beta_1& 1& 0&-1\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\beta_1\\in \\mathbf{F},\\]\n\t\\[A_{8}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1 & 0 & 0 & 0 \\\\\n\t0 & 1-\\alpha_1 & -\\alpha_1&0\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F},\\]\n\t\\[A_{9}=\\left(\n\t\\begin{array}{cccc}\n\t\\frac{1}{3}& 0 & 0 & 0 \\\\\n\t1 & \\frac{2}{3} & -\\frac{1}{3}&0\n\t\\end{array}\\right),\\\n\t\\ A_{10}=\\left(\n\t\\begin{array}{cccc}\n\t0 & 1 & 1 & 0 \\\\\n\t0 &0& 0 &-1\n\t\\end{array}\n\t\\right),\\]\n\t\\[ A_{11}=\\left(\n\t\\begin{array}{cccc}\n\t0 & 1 & 1 & 0 \\\\\n\t1 &0& 0 &-1\n\t\\end{array}\n\t\\right),\\\n\t\\ A_{12}=\\left(\n\t\\begin{array}{cccc}\n\t0 & 0 & 0 & 0 \\\\\n\t1 &0&0 &0\\end{array}\n\t\\right).\\]}\n\n{\\bf Proof.}\n{\\bf The first subset case.} In this case let $P(A)$ stand for a nonsingular matrix $\\left(\\begin{array}{cc} \\alpha_1+\\beta_3& \\alpha_2+\\beta_4\\\\ \\alpha_1+\\beta_2 & \\alpha_3+\\beta_4\\end{array}\\right)$ with rows\n\\[Tr_1(A)=(\\alpha _1+\\beta _3, \\alpha _2+\\beta _4),\\ Tr_2(A)=(\\alpha _1+\\beta _2, \\alpha _3+\\beta _4).\\] Due to (\\ref{2.1}) we have the identity \\[P(g A(g^{-1})^{\\otimes 2})= P(A)g^{-1},\\] which means that one can apply the Theorem as $V=\\mathbf{F}^8$ and $\\tau(g,A)= g A(g^{-1})^{\\otimes 2}.$ In this case for $V_0$ one can take \\[ V_0=\\{A\\in Mat(2\\times 4,\\mathbf{F}): \\det(P(A))=\\alpha_1(\\alpha_3-\\alpha_2)+\\beta_4(\\beta_3-\\beta_2)+\\alpha_3\\beta_3-\\alpha_2\\beta_2\\neq 0.\\}\\] Therefore,\ntwo-dimensional algebras $\\mathbf{A}$, $\\mathbf{B}$, given by\ntheir matrices of structure constants $A$, $B \\in V_0$, are isomorphic if and only if the equality \\[P(B)B(P(B)^{-1}\\otimes P(B)^{-1})=P(A)A(P(A)^{-1}\\otimes P(A)^{-1})\\] holds true.\n\nFor $\\xi_1=\\frac{\\alpha_3+\\beta_4}{\\Delta}$, $\\xi_2=-\\frac{\\alpha_1+\\beta_2}{\\Delta}$, $\\eta_1=-\\frac{\\alpha_2+\\beta_4}{\\Delta}$, $\\eta_2=\\frac{\\alpha_1+\\beta_3}{\\Delta}$, where $\\Delta=\\Delta(A)=\\det(P(A))$, one has\n\\[P(A)^{-1}=\\left(\\begin{array}{cc} \\xi_1& \\eta_1\\\\ \\xi_2& \\eta_2\\end{array}\\right)\\] and\n$A'=\\left(\\begin{array}{cccc} \\alpha'_1 & \\alpha'_2 & \\alpha'_3 &\\alpha'_4\\\\ \\beta'_1 & \\beta'_2 & \\beta'_3 &\\beta'_4\\end{array}\\right)=P(A)A(P(A)^{-1}\\otimes P(A)^{-1})$ is matrix consisting of columns\n\\[\\left(\\begin{array}{c}\n-\\frac{\\beta _1 \\eta _1 \\xi _1^2}{\\Delta }+\\frac{\\alpha _1 \\eta _2 \\xi _1^2}{\\Delta }+\\frac{2 \\alpha _1 \\eta _1 \\xi _1 \\xi _2}{\\Delta }-\\frac{2 \\beta _4 \\eta _2 \\xi _1 \\xi _2}{\\Delta }-\\frac{2 \\eta _1 \\eta _2 \\xi _1 \\xi _2}{\\Delta ^2}+\\frac{\\eta _2 \\xi _1^2 \\xi _2}{\\Delta ^2}-\\frac{\\beta _4 \\eta _1 \\xi _2^2}{\\Delta }+\\frac{\\alpha _4 \\eta _2 \\xi _2^2}{\\Delta }+\\frac{\\eta _1 \\xi _1 \\xi _2^2}{\\Delta ^2} \\\\\n\\frac{\\beta _1 \\xi _1^3}{\\Delta }-\\frac{3 \\alpha _1 \\xi _1^2 \\xi _2}{\\Delta }+\\frac{\\eta _2 \\xi _1^2 \\xi _2}{\\Delta ^2}+\\frac{3 \\beta _4 \\xi _1 \\xi _2^2}{\\Delta }+\\frac{\\eta _1 \\xi _1 \\xi _2^2}{\\Delta ^2}-\\frac{2 \\xi _1^2 \\xi _2^2}{\\Delta ^2}-\\frac{\\alpha _4 \\xi _2^3}{\\Delta }\n\\end{array}\n\\right),\\]\n\\[\\left(\\begin{array}{c}-\\frac{\\beta _1 \\eta _1^2 \\xi _1}{\\Delta }+\\frac{2 \\alpha _1 \\eta _1 \\eta _2 \\xi _1}{\\Delta }-\\frac{\\beta _4 \\eta _2^2 \\xi _1}{\\Delta }-\\frac{\\eta _1 \\eta _2^2 \\xi _1}{\\Delta ^2}+\\frac{\\alpha _1 \\eta _1^2 \\xi _2}{\\Delta }-\\frac{2 \\beta _4 \\eta _1 \\eta _2 \\xi _2}{\\Delta }-\\frac{\\eta _1^2 \\eta _2 \\xi _2}{\\Delta ^2}+\\frac{\\alpha _4 \\eta _2^2 \\xi _2}{\\Delta }+\\frac{2 \\eta _1 \\eta _2 \\xi _1 \\xi _2}{\\Delta ^2}\\\\\n\\frac{\\beta _1 \\eta _1 \\xi _1^2}{\\Delta }-\\frac{\\alpha _1 \\eta _2 \\xi _1^2}{\\Delta }-\\frac{2 \\alpha _1 \\eta _1 \\xi _1 \\xi _2}{\\Delta }+\\frac{2 \\beta _4 \\eta _2 \\xi _1 \\xi _2}{\\Delta }+\\frac{2 \\eta _1 \\eta _2 \\xi _1 \\xi _2}{\\Delta ^2}-\\frac{\\eta _2 \\xi _1^2 \\xi _2}{\\Delta ^2}+\\frac{\\beta _4 \\eta _1 \\xi _2^2}{\\Delta }-\\frac{\\alpha _4 \\eta _2 \\xi _2^2}{\\Delta }-\\frac{\\eta _1 \\xi _1 \\xi _2^2}{\\Delta ^2}\n\\end{array}\n\\right),\\]\n\n\\[\\left(\\begin{array}{c} -\\frac{\\beta _1 \\eta _1^2 \\xi _1}{\\Delta }+\\frac{2 \\alpha _1 \\eta _1 \\eta _2 \\xi _1}{\\Delta }-\\frac{\\beta _4 \\eta _2^2 \\xi _1}{\\Delta }-\\frac{\\eta _1 \\eta _2^2 \\xi _1}{\\Delta ^2}+\\frac{\\eta _2^2 \\xi _1^2}{\\Delta ^2}+\\frac{\\alpha _1 \\eta _1^2 \\xi _2}{\\Delta }-\\frac{2 \\beta _4 \\eta _1 \\eta _2 \\xi _2}{\\Delta }-\\frac{\\eta _1^2 \\eta _2 \\xi _2}{\\Delta ^2}+\\frac{\\alpha _4 \\eta _2^2 \\xi _2}{\\Delta }+\\frac{\\eta _1^2 \\xi _2^2}{\\Delta ^2}\\\\\n\\frac{\\beta _1 \\eta _1 \\xi _1^2}{\\Delta }-\\frac{\\alpha _1 \\eta _2 \\xi _1^2}{\\Delta }+\\frac{\\eta _2^2 \\xi _1^2}{\\Delta ^2}-\\frac{2 \\alpha _1 \\eta _1 \\xi _1 \\xi _2}{\\Delta }+\\frac{2 \\beta _4 \\eta _2 \\xi _1 \\xi _2}{\\Delta }-\\frac{\\eta _2 \\xi _1^2 \\xi _2}{\\Delta ^2}+\\frac{\\beta _4 \\eta _1 \\xi _2^2}{\\Delta }+\\frac{\\eta _1^2 \\xi _2^2}{\\Delta ^2}-\\frac{\\alpha _4 \\eta _2 \\xi _2^2}{\\Delta }-\\frac{\\eta _1 \\xi _1 \\xi _2^2}{\\Delta ^2}\\end{array}\n\\right),\\]\n\n\n\\[\\left(\\begin{array}{c} -\\frac{\\beta _1 \\eta _1^3}{\\Delta }+\\frac{3 \\alpha _1 \\eta _1^2 \\eta _2}{\\Delta }-\\frac{3 \\beta _4 \\eta _1 \\eta _2^2}{\\Delta }-\\frac{2 \\eta _1^2 \\eta _2^2}{\\Delta ^2}+\\frac{\\alpha _4 \\eta _2^3}{\\Delta }+\\frac{\\eta _1 \\eta _2^2 \\xi _1}{\\Delta ^2}+\\frac{\\eta _1^2 \\eta _2 \\xi _2}{\\Delta ^2}\\\\\n\\frac{\\beta _1 \\eta _1^2 \\xi _1}{\\Delta }-\\frac{2 \\alpha _1 \\eta _1 \\eta _2 \\xi _1}{\\Delta }+\\frac{\\beta _4 \\eta _2^2 \\xi _1}{\\Delta }+\\frac{\\eta _1 \\eta _2^2 \\xi _1}{\\Delta ^2}-\\frac{\\alpha _1 \\eta _1^2 \\xi _2}{\\Delta }+\\frac{2 \\beta _4 \\eta _1 \\eta _2 \\xi _2}{\\Delta }+\\frac{\\eta _1^2 \\eta _2 \\xi _2}{\\Delta ^2}-\\frac{\\alpha _4 \\eta _2^2 \\xi _2}{\\Delta }-\\frac{2 \\eta _1 \\eta _2 \\xi _1 \\xi _2}{\\Delta ^2}\\end{array}\n\\right).\\]\n\nFrom the above it can be easily seen that $\\alpha'_1=-\\beta'_2$, $\\alpha'_3=\\alpha'_2+1$, $\\beta'_3=\\beta'_2+1$ and $\\beta'_4=-\\alpha'_2,$ i.e.\n\\[A'=\\left(\n\\begin{array}{cccc}\n\\alpha'_1 & \\alpha'_2 &\\alpha'_2+1 & \\alpha'_4 \\\\\n\\beta'_1 & -\\alpha'_1 & -\\alpha'_1+1 & -\\alpha'_2\n\\end{array}\\right).\\]\nTherefore, the main role in the finding of $A'$ is played by the functions\n\n$\\alpha'_1=-\\frac{\\beta _1 \\eta _1 \\xi _1^2}{\\Delta }+\\frac{\\alpha _1 \\eta _2 \\xi _1^2}{\\Delta }+\\frac{2 \\alpha _1 \\eta _1 \\xi _1 \\xi _2}{\\Delta }-\\frac{2 \\beta _4 \\eta _2 \\xi _1 \\xi _2}{\\Delta }-\\frac{2 \\eta _1 \\eta _2 \\xi _1 \\xi _2}{\\Delta ^2}+\\frac{\\eta _2 \\xi _1^2 \\xi _2}{\\Delta ^2}-\\frac{\\beta _4 \\eta _1 \\xi _2^2}{\\Delta }+\\frac{\\alpha _4 \\eta _2 \\xi _2^2}{\\Delta }+\\frac{\\eta _1 \\xi _1 \\xi _2^2}{\\Delta ^2},$\n\n$\\beta'_1=\\frac{\\beta _1 \\xi _1^3}{\\Delta }-\\frac{3 \\alpha _1 \\xi _1^2 \\xi _2}{\\Delta }+\\frac{\\eta _2 \\xi _1^2 \\xi _2}{\\Delta ^2}+\\frac{3 \\beta _4 \\xi _1 \\xi _2^2}{\\Delta }+\\frac{\\eta _1 \\xi _1 \\xi _2^2}{\\Delta ^2}-\\frac{2 \\xi _1^2 \\xi _2^2}{\\Delta ^2}-\\frac{\\alpha _4 \\xi _2^3}{\\Delta }\n,$\n\n$\\alpha'_2=-\\frac{\\beta _1 \\eta _1^2 \\xi _1}{\\Delta }+\\frac{2 \\alpha _1 \\eta _1 \\eta _2 \\xi _1}{\\Delta }-\\frac{\\beta _4 \\eta _2^2 \\xi _1}{\\Delta }-\\frac{\\eta _1 \\eta _2^2 \\xi _1}{\\Delta ^2}+\\frac{\\alpha _1 \\eta _1^2 \\xi _2}{\\Delta }-\\frac{2 \\beta _4 \\eta _1 \\eta _2 \\xi _2}{\\Delta }-\\frac{\\eta _1^2 \\eta _2 \\xi _2}{\\Delta ^2}+\\frac{\\alpha _4 \\eta _2^2 \\xi _2}{\\Delta }+\\frac{2 \\eta _1 \\eta _2 \\xi _1 \\xi _2}{\\Delta ^2},$\n\n$\\alpha'_4=-\\frac{\\beta _1 \\eta _1^3}{\\Delta }+\\frac{3 \\alpha _1 \\eta _1^2 \\eta _2}{\\Delta }-\\frac{3 \\beta _4 \\eta _1 \\eta _2^2}{\\Delta }-\\frac{2 \\eta _1^2 \\eta _2^2}{\\Delta ^2}+\\frac{\\alpha _4 \\eta _2^3}{\\Delta }+\\frac{\\eta _1 \\eta _2^2 \\xi _1}{\\Delta ^2}+\\frac{\\eta _1^2 \\eta _2 \\xi _2}{\\Delta ^2}.$\n\nIt can be easily verified that these functions can take any values in $\\mathbf{F}.$ Therefore, the values of $\\alpha'_1,\\beta'_1, \\alpha'_2, \\alpha'_4$ also can be any elements in $\\mathbf{F}$. Using obvious redenotions one can list the following non-isomorphic ``canonical'' algebras from the first subset, given by their MSCs as follows:\n\\[ A_{1}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & \\alpha_2 &\\alpha_2+1 & \\alpha_4 \\\\\n\\beta_1 & -\\alpha_1 & -\\alpha_1+1 & -\\alpha_2\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\alpha_2, \\alpha_4, \\beta_1) \\in \\mathbf{F}^4.\\]\n\nFor any algebra $\\mathbf{A}$ from the first subset there exists a unique algebra from the class $A_{1}(\\mathbf{c})$ isomorphic to $\\mathbf{A}$. Note that the appearance of $A_{1}(\\mathbf{c})$ for the first subset does not depend on either algebraically closeness of $\\mathbf{F}$ or its characteristic.\n\n{\\bf The second and third subset cases.} In these cases one can make $Tr_1(A)g=(1,0)$. It implies that $Tr_2(A)g=(\\lambda,0)$. Therefore, we can assume that\n\\[Tr_1(A)=(\\alpha _1+\\beta _3,\\alpha _2+\\beta _4)=(1,0)\\ \\mbox{and}\\ Tr_2(A)=(\\lambda(\\alpha _1+\\beta _3),\\lambda(\\alpha _2+\\beta _4))=(\\lambda,0).\\]\nHere if $\\lambda$ is zero we can also cover the third subset. So let us consider\n\n\\[A=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & \\alpha _2 & \\lambda \\alpha _2+(\\lambda -1)\\beta _4 & \\alpha _4 \\\\\n\\beta _1 & (\\lambda -1)\\alpha _1+\\lambda \\beta _3 & \\beta _3 & \\beta _4\n\\end{array}\n\\right)=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & \\alpha _2 & \\alpha _2 & \\alpha _4 \\\\\n\\beta _1 & \\lambda-\\alpha _1 & 1-\\alpha _1& -\\alpha _2\n\\end{array}\\right),\\] with respect to $g\\in GL(2,\\mathbf{F})$ of the form \\[g^{-1}=\\left(\\begin{array}{cc} 1& 0\\\\ \\xi_2& \\eta_2\\end{array}\\right),\\] as far as\n$(1,0)g^{-1}=(1,0)$ if and only if $g^{-1}$ is of the above form. In this case for the entries of $A'$ we have $A'=\\left(\n\\begin{array}{cccc}\n\\alpha'_1 & \\alpha'_2 & \\alpha'_2 & \\alpha'_4 \\\\\n\\beta'_1 & \\lambda-\\alpha'_1 & 1-\\alpha'_1& -\\alpha'_2\n\\end{array}\\right)=gA(g^{-1})^{\\otimes 2}$, where $\\lambda$ is same for the $A$ and $A'$,\none has\n\n$\\alpha '_1=\\frac{1}{\\Delta }(\\alpha _1 \\eta _2 +2\\alpha _2\\eta _2\\xi _2+\\alpha _4 \\eta _2 \\xi _2^2)=\\alpha _1+2\\alpha _2\\xi _2+\\alpha _4\\xi _2^2.$\n\n$\\alpha '_2=\\frac{-1}{\\Delta }(-\\alpha _2\\eta^2_2 -\\alpha _4 \\eta _2^2 \\xi _2)=(\\alpha _2+\\alpha _4\\xi _2)\\eta _2.$\n\n$\\alpha '_4=\\frac{-1}{\\Delta }(\\beta _1 \\eta _1^3+[(\\lambda-2)\\alpha _1+(1+\\lambda)\\beta _3]\\eta _1^2 \\eta _2-[(1+\\lambda)\\alpha _2+(\\lambda-2)\\beta _4]\\eta _1 \\eta _2^2-\\alpha _4 \\eta _2^3)=\\alpha _4 \\eta_2^2.$\n\n$\\beta '_1=\n\\frac{1}{\\Delta }(\\beta _1 \\xi _1^3+[(\\lambda-2)\\alpha _1 +(\\lambda+1)\\beta _3]\\xi _1^2 \\xi _2-[(1+\\lambda)\\alpha _2+(\\lambda-2)\\beta _4] \\xi _1 \\xi _2^2-\\alpha _4 \\xi _2^3)$\n\n\\qquad $=\\frac{\\beta _1+(1+\\lambda-3 \\alpha _1)\\xi _2-3 \\alpha _2 \\xi _2^2-\\alpha _4 \\xi _2^3}{\\eta_2}.$\n\nSo it is enough to consider the system\n\n$\\alpha '_1=\\alpha _1+2 \\alpha _2 \\xi _2+\\alpha _4 \\xi _2^2,$\n\n$\\alpha '_2=(\\alpha _2+\\alpha_4\\xi_2) \\eta _2,$\n\n$\\alpha '_4=\\alpha _4 \\eta _2^2,$\n\n$\\beta '_1=\\frac{\\beta _1+(1+\\lambda-3\\alpha _1)\\xi _2-3 \\alpha _2 \\xi _2^2-\\alpha _4 \\xi _2^3}{\\eta _2}.$\n\nNow we have to consider a few cases again.\n\nCase 1: $\\alpha_4\\neq 0$. In this case one can make $$\\alpha '_2=0, \\alpha '_4=1, \\alpha '_1=\\alpha_1-\\frac{\\alpha^2_2}{\\alpha_4}\\ \\mbox{and}\\ \\beta '_1=\\sqrt{\\alpha_4}\\left(\\beta _1-\\left(1+\\lambda-3\\alpha _1\\right)\\frac{\\alpha_2}{\\alpha_4}-2\\frac{\\alpha^3_2}{\\alpha^2_4}\\right)$$\nand\nonce again using redenotion one can represent the corresponding ``canonical'' MSCs as follows\n\\[ A_{2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n\\beta _1& \\beta_2 & 1-\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1, \\beta_2)\\in \\mathbf{F}^3,\\] as far as the expressions $\\alpha_1-\\frac{\\alpha^2_2}{\\alpha_4}$, $\\sqrt{\\alpha_4}(\\beta _1-(1+\\lambda-3\\alpha _1)\\frac{\\alpha_2}{\\alpha_4}-2\\frac{\\alpha^3_2}{\\alpha^2_4})$, $\\lambda-(\\alpha_1-\\frac{\\alpha^2_2}{\\alpha_4} ) $ may have any values in $\\mathbf{F}$. It can be checked that algebras $\\left(\n\\begin{array}{cccc}\n\t\\alpha_1 & 0 & 0 & 1 \\\\\n\t\\beta _1& \\beta_2 & 1-\\alpha_1&0\n\\end{array}\\right)$ and $\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n-\\beta _1& \\beta_2 & 1-\\alpha_1&0\n\\end{array}\\right)$ are isomorphic.\nNote that the result $A_{2}(\\mathbf{c})$, does not depend on the characteristic of $\\mathbf{F}$.\n\nCase 2: $\\alpha _4=0.$\n\nSubcase 2 - a): $\\alpha _2\\neq 0.$ If $\\alpha _2\\neq 0$ then one can make\n$\\alpha '_1=0$, $\\alpha '_2=1$, $\\beta '_1=\\alpha_2\\beta _1-\\frac{2+2\\lambda-3\\alpha _1}{4}\\alpha_1$ to get the following set of canonical matrices of structural constants\n\\[ A_{3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n\\beta _1& \\beta_2 & 1&-1\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\beta_1, \\beta_2)\\in \\mathbf{F}^2.\\]\n\nSubcase 2 - b): $\\alpha _2= 0.$ If $\\alpha _2=0$ then $\\alpha '_1=\\alpha _1,$ $\\alpha '_2=0,$ $\\alpha '_4=0,$ $\\beta '_1=\\frac{\\beta _1+(1+\\lambda-3\\alpha _1)\\xi _2}{\\eta _2}.$\n\n\\underline{Subsubcase: 2 - b) - 1: $1+\\lambda-3\\alpha _1\\neq 0.$} If $1+\\lambda-3\\alpha _1\\neq 0$, that is $\\lambda-\\alpha _1\\neq 2\\alpha_1-1$, one can make $\\beta'_1=0$ to get\n\\[A_{4}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & 0 & 0 & 0 \\\\\n0 & \\beta_2 & 1-\\alpha _1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_2) \\in \\mathbf{F}^2, \\ \\mbox{with}\\ \\beta_2\\neq 2\\alpha_1-1.\\]\n\n\\underline{Subsubcase: 2 - b) - 2: $1+\\lambda-3\\alpha _1= 0.$} If $1+\\lambda-3\\alpha _1=0$ and $\\beta_1\\neq 0$ one can make $\\beta'_1=1$ to get\n\\[A_{5}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & 0 & 0 & 0 \\\\\n1 & 2\\alpha _1-1 & 1-\\alpha _1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in F.\\] If $\\beta_1=0$ one has $\\lambda-\\alpha _1= 2\\alpha_1-1$ and therefore\n$A'=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & 0 & 0 & 0 \\\\\n0 & 2\\alpha _1-1 & 1-\\alpha _1&0\n\\end{array}\\right)$\\\\ which is $A_{4}(\\mathbf{c})$ with $\\beta_2=2\\alpha_1-1$.\n\n{\\bf The fourth subset case.}\nBy the similar justification as in the second and the third subsets case it is enough to consider \\[ A=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & \\alpha _2 & \\alpha _2 & \\alpha _4 \\\\\n\\beta _1 & 1-\\alpha _1 & -\\alpha _1 & -\\alpha _2\n\\end{array}\n\\right)\\ \\mbox{and}\\ \\ g^{-1}=\\left(\\begin{array}{cc} 1& 0\\\\ \\xi_2& \\eta_2\\end{array}\\right),\\\n\\mbox{where one has}\\]\n\n$\\alpha'_1= \\alpha _1+2 \\alpha _2 \\xi _2+\\alpha _4 \\xi _2^2 ,$\n\n$\\alpha'_2= (\\alpha _2+\\alpha _4\\xi _2)\\eta_2,$\n\n$\\alpha'_4=\\alpha _4 \\eta _2^2,$\n\n$\\beta'_1=\\frac{\\beta _1+\\xi _2-3 \\alpha _1 \\xi _2-3 \\alpha _2 \\xi _2^2-\\alpha _4 \\xi _2^3}{\\eta _2}.$\n\nTherefore we get the canonical MSCs as follows\n\\[A_{6}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right)\\simeq \\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n-\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1) \\in \\mathbf{F}^2,\\]\n(it is $A_{2}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\alpha_1+\\beta_2=0$) or\n\\[ A_{7}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n\\beta _1& 1& 0&-1\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\beta_1\\in \\mathbf{F},\\]\n(it is $A_{3}(\\mathbf{\\mathbf{c}})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\beta_2=0$) or\n\\[A_{8}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & 0 & 0 & 0 \\\\\n0 & 1-\\alpha _1 & -\\alpha _1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F},\\]\n(it is $A_{4}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\alpha_1+\\beta_2=0$) or\n\\[A_{9}=\\left(\n\\begin{array}{cccc}\n\\frac{1}{3} & 0 & 0 & 0 \\\\\n1 & \\frac{2}{3} & -\\frac{1}{3} &0\n\\end{array}\\right)\\] (it is $A_{5}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $3\\alpha_1-1=0$).\n\n{\\bf The fifth subset case.} In this case\n\\[A=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & \\alpha _2 & \\alpha _2 & \\alpha _4 \\\\\n\\beta _1 & -\\alpha _1 & -\\alpha _1 & -\\alpha _2\n\\end{array}\n\\right),\\ g^{-1}=\\left(\\begin{array}{cc} \\xi_1& \\eta_1\\\\ \\xi_2& \\eta_2\\end{array}\\right)\\mbox{and}\\\nA'=\\left(\n\\begin{array}{cccc}\n\\alpha' _1 & \\alpha' _2 & \\alpha' _2 & \\alpha' _4 \\\\\n\\beta' _1 & -\\alpha' _1 & -\\alpha' _1 & -\\alpha' _2\n\\end{array}\n\\right)\\]\nsuch that:\n\n$\\alpha' _1=\\frac{1}{\\Delta }\\left(-\\beta _1 \\eta _1 \\xi _1^2+\\alpha _1 \\eta _2 \\xi _1^2+2 \\alpha _1 \\eta _1 \\xi _1 \\xi _2+2 \\alpha _2 \\eta _2 \\xi _1 \\xi _2+\\alpha _2 \\eta _1 \\xi _2^2+\\alpha _4 \\eta _2 \\xi _2^2\\right),$\n\n$\\alpha' _2=\\frac{-1}{\\Delta }\\left(\\beta _1 \\eta _1^2 \\xi _1-2 \\alpha _1 \\eta _1 \\eta _2 \\xi _1-\\alpha _2 \\eta _2^2 \\xi _1-\\alpha _1 \\eta _1^2 \\xi _2-2 \\alpha _2 \\eta _1 \\eta _2 \\xi _2-\\alpha _4 \\eta _2^2 \\xi _2\\right),$\n\n$\\alpha' _4=\\frac{-1}{\\Delta }\\left(\\beta _1 \\eta _1^3-3 \\alpha _1 \\eta _1^2 \\eta _2-3 \\alpha _2 \\eta _1 \\eta _2^2-\\alpha _4 \\eta _2^3\\right),$\n\n$\\beta' _1=\\frac{1}{\\Delta }\\left(\\beta _1 \\xi _1^3-3 \\alpha _1 \\xi _1^2 \\xi _2-3 \\alpha _2 \\xi _1 \\xi _2^2-\\alpha _4 \\xi _2^3\\right).$\n\nIf $\\alpha _4\\neq 0$ by making $\\frac{\\eta_2}{\\eta_1}$ equal to any root of the polynomial $p(t)=\\beta _1-3 \\alpha _1t-3 \\alpha _2t^2-\\alpha _4t^3$ one can make $\\alpha'_4=0$. Therefore, further it is assumed that $\\alpha _4= 0$.\n\nLet us consider $g$ with $\\eta _1 = 0$ to have $\\alpha'_4=0.$ In this case $\\Delta=\\xi_1\\eta_2$ and \\\\\n$\\alpha' _1=\\xi _1\\left(\\alpha _1+2\\alpha _2\\frac{\\xi_2}{\\xi_1}\\right),$\\\\\n$\\alpha' _2=\\alpha _2\\eta_2,$\\\\\n$\\beta' _1=\\frac{ \\xi _1^2}{\\eta_2 }\\left(\\beta _1-3 \\alpha _1 \\frac{\\xi_2}{\\xi_1}-3 \\alpha _2(\\frac{\\xi_2}{\\xi_1})^2\\right).$\\\\\nCase a: $\\alpha_2\\neq 0$. One can consider $\\frac{\\xi_2}{\\xi_1}=\\frac{-\\alpha _1}{2\\alpha _2}$ to get $\\alpha' _1=0, \\alpha' _2=1$ and $\\beta' _1=\\xi^2_1\\frac{3\\alpha^2 _1+4\\alpha _2\\beta _1}{4}.$ Therefore one can make $\\beta' _1$ equal to $0$ or $1,$ depending on $3\\alpha^2 _1+4\\alpha _2\\beta _1$ to have\n\\[A_{10}=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n0 &0& 0 &-1\n\\end{array}\n\\right)\\ \\mbox{or}\\ A_{11}=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n1 &0& 0 &-1\n\\end{array}\n\\right).\\]\nCase b: $\\alpha_2=0$. Then $\\alpha'_2=\\alpha'_4=0$ and $\\alpha'_1=\\xi_1\\alpha_1, \\beta' _1=\\frac{\\xi^2_1}{\\eta_2}\\left(\\beta _1-3 \\alpha _1 \\frac{\\xi_2}{\\xi_1}\\right).$ Therefore if $\\alpha_1\\neq 0$ one can make $\\alpha'_1=1, \\beta' _1=0$ to get $A'=\\left(\n\\begin{array}{cccc}\n1 & 0 & 0 & 0 \\\\\n0 &-1& -1 &0\n\\end{array}\n\\right)$ , which isomorphic to $A_{10},$ if $\\alpha_1= 0$ then $\\alpha'_1=0$ and one can make $\\beta' _1=1$ to come to\n\\[A_{12}=\\left(\n\\begin{array}{cccc}\n0 & 0 & 0 & 0 \\\\\n1 &0&0 &0\\end{array}\n\\right)\\]\nA routine check in each subset case shows that the corresponding algebras presented above\nare not isomorphic. \n\n\\subsection{Characteristic $2$ case}\nIn this case the result can be summarized in the following form.\n\n{\\bf Theorem 2.}\\textit{ Over an algebraically closed field $\\mathbf{F}$ characteristic 2 any non-trivial 2-dimensional algebra is isomorphic to only one of the following listed, by their matrices of structure constants, algebras:\n\\[A_{1,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & \\alpha_2 &\\alpha_2+1 & \\alpha_4 \\\\\n\\beta_1 & -\\alpha_1 & -\\alpha_1+1 & -\\alpha_2\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\alpha_2, \\alpha_4, \\beta_1)\\in \\mathbf{F}^4,\\]\n\\[A_{2,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n\\beta _1& \\beta_2 & 1-\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1, \\beta_2)\\in \\mathbf{F}^3,\\]\n\\[A_{3,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 1 & 1 & 0 \\\\\n0& \\beta_2 & 1-\\alpha_1&1\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_2)\\in \\mathbf{F}^2,\\]\n\\[A_{4,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & 0 & 0 & 0 \\\\\n0 & \\beta_2 & 1-\\alpha _1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1,\\beta_2)\\in \\mathbf{F}^2,\\]\n\\[A_{5,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 0 \\\\\n1 & 1 & 1-\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F},\\]\n\\[A_{6,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1)\\in \\mathbf{F}^2,\\]\n\\[A_{7,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 1 & 1 & 0 \\\\\n0& 1-\\alpha_1& -\\alpha_1&-1\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F},\\]\n\\[A_{8,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 0 \\\\\n0 & 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F},\\]\n\\[A_{9,2}=\\left(\n\\begin{array}{cccc}\n1 & 0 & 0 & 0 \\\\\n1 & 0 & 1&0\n\\end{array}\\right),\\ \\ A_{10,2}=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n0 &0& 0 &-1\n\\end{array}\n\\right),\\]\n\\[ A_{11,2}=\\left(\n\\begin{array}{cccc}\n1 & 1 & 1 & 0 \\\\\n0 &-1& -1 &-1\n\\end{array}\n\\right),\\\n\\ A_{12,2}=\\left(\n\\begin{array}{cccc}\n0 & 0 & 0 & 0 \\\\\n1 &0&0 &0\\end{array}\n\\right).\\]}\n\n{\\bf Proof.}\nIn this case $\\alpha=-\\alpha$ for any $\\alpha\\in \\mathbf{F}$, but sometimes we use $-\\alpha$ also to keep similarity with the previous case.\n\n{\\bf The first subset case.} It has been noted earlier that in this case the result doesn't depend either on algebraically closeness of $\\mathbf{F}$ or its characteristics. Therefore, one has \\[A_{1,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & \\alpha_2 &\\alpha_2+1 & \\alpha_4 \\\\\n\\beta_1 & -\\alpha_1 & -\\alpha_1+1 & -\\alpha_2\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\alpha_2, \\alpha_4, \\beta_1)\\in \\mathbf{F}^4.\\]\n\n{\\bf The second and third subset cases.}\nObviously in this case one can make $Tr_1(A)g=(1,0)$. This implies $Tr_2(A)g=(\\lambda,0)$. Therefore, it can be assumed that $Tr_1(A)=(\\alpha _1+\\beta _3,\\alpha _2+\\beta _4)=(1,0)$ and $Tr_2(A)=( \\lambda(\\alpha _1+\\beta _3),\\lambda(\\alpha _2+\\beta _4)=(\\lambda,0)$. Here also we allow $\\lambda$ to be zero to include the third subset's case. Therefore, we consider\n\\[A=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & \\alpha _2 & \\alpha_2 & \\alpha _4 \\\\\n\\beta _1 & \\lambda-\\alpha _1 & 1-\\alpha _1& -\\alpha _2\n\\end{array}\\right)\\ \\mbox{with respect to}\\ g^{-1}=\\left(\\begin{array}{cc} 1& 0\\\\ \\xi_2& \\eta_2\\end{array}\\right) \\ \\mbox{and}\\]\n\n$\\alpha '_1=\\alpha _1+\\alpha _4\\xi _2^2,$\n\n$\\alpha '_2=(\\alpha _2+\\alpha _4\\xi _2)\\eta _2,$\n\n$\\alpha '_4=\\alpha _4 \\eta_2^2,$\n\n$\\beta '_1=\\frac{\\beta _1+(1+\\lambda-3 \\alpha _1)\\xi _2-3 \\alpha _2 \\xi _2^2-\\alpha _4 \\xi _2^3}{\\eta_2}=\\frac{\\beta _1+(1+\\lambda+\\alpha _1)\\xi _2+ \\alpha _2 \\xi _2^2+\\alpha _4 \\xi _2^3}{\\eta_2}.$\n\\begin{itemize}\n\t\\item[1.] As we have observed earlier in $\\alpha_4\\neq 0$ case one gets\n\t\\[A_{2,2}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1 & 0 & 0 & 1 \\\\\n\t\\beta _1& \\beta_2 & 1-\\alpha_1&0\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1, \\beta_2)\\in \\mathbf{F}^3.\\]\n\t\\item[2.] Let now assume that $\\alpha _4=0.$\n\t\\begin{itemize}\n\t\t\\item[2-a).] If $\\alpha _2\\neq0$ then one can make\n\t\t$\\alpha '_1=\\alpha_1$, $\\alpha '_2=1$, $\\alpha'_4=0$ and $\\beta'_1=0$ to get \\[A_{3,2}(\\mathbf{c})=\\left(\n\t\t\\begin{array}{cccc}\n\t\t\\alpha_1 & 1 & 1 & 0 \\\\\n\t\t0& \\beta_2 & 1-\\alpha_1&-1\n\t\t\\end{array}\\right), \\ \\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_2)\\in \\mathbf{F}^2.\\]\n\t\t\n\t\t\\item[2-b).] If $\\alpha _2=0$ then $\\alpha '_1=\\alpha _1,$ $\\alpha '_2=0,$ $\\alpha '_4=0,$ and $\\beta '_1=\\frac{\\beta _1+(1+\\lambda+\\alpha _1)\\xi _2}{\\eta _2}.$\n\t\t\\begin{itemize}\n\t\t\t\\item[2-b)-1.] If $1+\\lambda+\\alpha _1\\neq 0$, that is $\\lambda-\\alpha _1\\neq 1$, one can make $\\beta'_1=0$ to get\n\t\t\t\\[A_{4,2}(\\mathbf{c})=\\left(\n\t\t\t\\begin{array}{cccc}\n\t\t\t\\alpha _1 & 0 & 0 & 0 \\\\\n\t\t\t0 & \\beta_2 & 1-\\alpha _1&0\n\t\t\t\\end{array}\\right),\\ \\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_2)\\in \\mathbf{F}\\ \\mbox{with}\\ \\beta_2\\neq 1.\\]\n\t\t\t\n\t\t\t\\item[2-b)-2.] Let $1+\\lambda+\\alpha _1=0$. In this case\n\t\t\t\\begin{itemize}\n\t\t\t\t\\item if $\\beta_1\\neq 0$ one can make $\\beta'_1=1$ to have\n\t\t\t\t\\[A_{5,2}(\\mathbf{c})=\\left(\n\t\t\t\t\\begin{array}{cccc}\n\t\t\t\t\\alpha _1 & 0 & 0 & 0 \\\\\n\t\t\t\t1 & 1 & 1-\\alpha _1&0\n\t\t\t\t\\end{array}\\right), \\ \\mbox{where}\\ \\mathbf{c}= \\alpha _1 \\in \\mathbf{F}.\\]\n\t\t\t\t\\item if $\\beta_1=0$ then\n\t\t\t\t\\[A'=\\left(\n\t\t\t\t\\begin{array}{cccc}\n\t\t\t\t\\alpha _1 & 0 & 0 & 0 \\\\\n\t\t\t\t0 & 1 & 1-\\alpha _1&0\n\t\t\t\t\\end{array}\\right), \\ \\mbox{which is}\\ A_{4,2}(\\mathbf{c}) \\ \\mbox{with}\\ \\beta_2=1.\\]\n\t\t\t\\end{itemize}\n\t\t\\end{itemize}\n\t\\end{itemize}\n\\end{itemize}\n{\\bf The fourth subset case.} For this case we have \\[ A=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & \\alpha _2 & \\alpha _2 & \\alpha _4 \\\\\n\\beta _1 & 1-\\alpha _1 & -\\alpha _1 & -\\alpha _2\n\\end{array}\n\\right)\\ \\mbox{and}\\ g^{-1}=\\left(\\begin{array}{cc} 1& 0\\\\ \\xi_2& \\eta_2\\end{array}\\right)\\]\ntherefore\n\n$\\alpha'_1= \\alpha _1+\\alpha _4 \\xi _2^2 ,$\n\n$\\alpha'_2= (\\alpha _2+\\alpha _4\\xi _2)\\eta_2,$\n\n$\\alpha'_4=\\alpha _4 \\eta _2^2,$\n\n$\\beta'_1=\\frac{\\beta _1+(1+\\alpha _1)\\xi _2+\\alpha _2 \\xi _2^2+\\alpha _4 \\xi _2^3}{\\eta _2}.$\n\nHence, one gets\n$A_{6,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right) \\ \\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1)\\in \\mathbf{F}^2$\\\\ (it is $A_{2,2}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\alpha_1+\\beta_2=0$),\n\\[A_{7,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 1 & 1 & 0 \\\\\n0&1-\\alpha_1&-\\alpha_1&-1\n\\end{array}\\right),\\ \\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F}\\] (it is $A_{3,2}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\alpha_1+\\beta_2=0$),\n\\[A_{8,2}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 0 \\\\\n0 & 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right)\\ \\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F}\\] (it is $A_{4,2}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\alpha_1+\\beta_2=0$)\nand\n\\[A_{9,2}=\\left(\n\\begin{array}{cccc}\n1 & 0 & 0 & 0 \\\\\n1 &0 & 1&0\n\\end{array}\\right)\\] (it is $A_{5,2}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\alpha_1+1=0$).\n\n{\\bf The fifth subset case.} In this case\n\\[A=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & \\alpha _2 & \\alpha _2 & \\alpha _4 \\\\\n\\beta _1 & -\\alpha _1 & -\\alpha _1 & -\\alpha _2\n\\end{array}\\right),\\ \\ g^{-1}=\\left(\\begin{array}{cc} \\xi_1& \\eta_1\\\\ \\xi_2& \\eta_2\\end{array}\\right)\\]\nand for the entries of $A'=\\left(\n\\begin{array}{cccc}\n\\alpha' _1 & \\alpha' _2 & \\alpha' _2 & \\alpha' _4 \\\\\n\\beta' _1 & -\\alpha' _1 & -\\alpha' _1 & -\\alpha' _2\n\\end{array}\n\\right)$ one has\n\n$\\alpha' _1=\\frac{1}{\\Delta }\\left(\\beta _1 \\eta _1 \\xi _1^2+\\alpha _1 \\eta _2 \\xi _1^2+ \\alpha_2 \\eta _1\\xi^2_2+\\alpha _4 \\eta _2 \\xi _2^2\\right),$\n\n$\\alpha' _2=\\frac{1}{\\Delta }\\left(\\beta _1 \\eta _1^2 \\xi _1+\\alpha _1 \\eta _1^2 \\xi _2+\\alpha _2 \\eta _2^2 \\xi _1+\\alpha _4 \\eta _2^2 \\xi _2\\right),$\n\n$\\alpha' _4=\\frac{1}{\\Delta }\\left(\\beta _1 \\eta _1^3+ \\alpha _1 \\eta _1^2 \\eta _2+\\alpha _2 \\eta _1 \\eta _2^2+\\alpha _4 \\eta _2^3\\right),$\n\n$\\beta' _1=\\frac{1}{\\Delta }\\left(\\beta _1 \\xi _1^3+\\alpha _1 \\xi _1^2 \\xi _2+\\alpha _2 \\xi _1 \\xi _2^2+\\alpha _4 \\xi _2^3\\right).$\n\nAs it was observed one can assume that $\\alpha_4=0$.\n\nLet us consider $g$ with $\\eta _1 = 0$ to have $\\alpha'_4=0.$ In this case $\\Delta=\\xi_1\\eta_2$ and \\\\\n$\\alpha' _1=\\xi _1\\alpha _1,$\\\\\n$\\alpha' _2=\\alpha _2\\eta_2,$\\\\\n$\\beta' _1=\\frac{ \\xi _1^2}{\\eta_2 }\\left(\\beta _1+ \\alpha _1 \\frac{\\xi_2}{\\xi_1}+ \\alpha _2(\\frac{\\xi_2}{\\xi_1})^2\\right).$\\\\\nIf one of $\\alpha _1, \\alpha _2$ is not zero one can make $\\beta' _1=0$ and depending on $\\alpha _1, \\alpha _2$ one can have $\\alpha' _1=0, \\alpha' _2=1$ or $\\alpha' _1=1, \\alpha' _2=1$ or $\\alpha' _1=1, \\alpha' _2=0$ to get\n$A_{10,2}=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n0 &0& 0 &-1\n\\end{array}\n\\right)$ or $A_{11,2}=\\left(\n\\begin{array}{cccc}\n1 & 1 & 1 & 0 \\\\\n0 &-1& -1 &-1\n\\end{array}\n\\right)$ or $A'=\\left(\n\\begin{array}{cccc}\n1 & 0 & 0 & 0 \\\\\n0 &-1& -1 &0\n\\end{array}\n\\right)$ which is isomorphic to $A_{10,2}$\\\\\nIf $\\alpha _1=\\alpha _2=0$ then one can make $\\beta' _1=1$ to get\n\\[A_{12,2}=\\left(\n\\begin{array}{cccc}\n0 & 0 & 0 & 0 \\\\\n1 &0& 0 &0\n\\end{array}\n\\right).\\]\n\n\\subsection{Characteristic $3$ case}\nIn this case we summarize the result as follows.\n\n{\\bf Theorem 3.}\\textit{ Over an algebraically closed field $\\mathbf{F}$ characteristics 3 any non-trivial 2-dimensional algebra is isomorphic to only one of the following listed, by their matrices of structure constant matrices, algebras:\n\\[A_{1,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & \\alpha_2 &\\alpha_2+1 & \\alpha_4 \\\\\n\\beta_1 & -\\alpha_1 & -\\alpha_1+1 & -\\alpha_2\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\alpha_2, \\alpha_4, \\beta_1)\\in \\mathbf{F}^4,\\]\n\\[A_{2,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n\\beta _1& \\beta _2& 1-\\alpha_1&0\n\\end{array}\\right)\\simeq\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n-\\beta _1& \\beta _2& 1-\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1, \\beta_2)\\in \\mathbf{F}^3,\\]\n\\[A_{3,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n\\beta _1& \\beta _2 & 1&-1\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\beta_1, \\beta_2)\\in \\mathbf{F}^2,\\]\n\\[A_{4,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & 0 & 0 & 0 \\\\\n0 & \\beta _2& 1-\\alpha _1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_2)\\in \\mathbf{F}^2,\\]\n\\[A_{5,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1& 0 & 0 & 0 \\\\\n1 & -1-\\alpha_1 & 1-\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F},\\]\n\\[A_{6,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right)\\simeq \\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n-\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1)\\in \\mathbf{F}^2,\\]\n\\[A_{7,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n\\beta_1& 1& 0&-1\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\beta_1\\in \\mathbf{F},\\]\n\\[A_{8,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 0 \\\\\n0 & 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F},\\]\n\\[A_{9,3}=\\left(\n\\begin{array}{cccc}\n0 & 1& 1& 0 \\\\\n1 &0&0 &-1\\end{array}\n\\right),\\\n\\ A_{10,3}=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n0 &0&0 &-1\\end{array}\n\\right),\\]\n\\[ A_{11,3}=\\left(\n\\begin{array}{cccc}\n1 & 0 & 0 & 0 \\\\\n1 &-1&-1 &0\\end{array}\n\\right),\\ \\ A_{12,3}=\\left(\n\\begin{array}{cccc}\n0 &0 &0 & 0 \\\\\n1 &0& 0 &0\n\\end{array}\n\\right).\\]}\n\n{\\bf Proof.}\n\\quad {\\bf The first subset case.} One has\n$$A_{1,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & \\alpha_2 &\\alpha_2+1 & \\alpha_4 \\\\\n\\beta_1 & -\\alpha_1 & -\\alpha_1+1 & -\\alpha_2\n\\end{array}\\right), \\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\alpha_2, \\alpha_4, \\beta_1)\\in \\mathbf{F}^4.$$\n\n{\\bf The second and third subset cases.} One can consider\n\\[A=\\left(\n\\begin{array}{cccc}\n\\alpha _1 & \\alpha _2 & \\alpha _2 & \\alpha _4 \\\\\n\\beta _1 & \\lambda-\\alpha _1 & 1-\\alpha _1& -\\alpha _2\n\\end{array}\\right),\\] where $\\lambda$ may be zero as well, with respect to\n$g^{-1}=\\left(\\begin{array}{cc} 1& 0\\\\ \\xi_2& \\eta_2\\end{array}\\right)$ in that case for the entries of $A'$ we get\n\n$\\alpha '_1=\\alpha _1+2 \\alpha _2 \\xi _2+\\alpha _4 \\xi _2^2,$\n\n$\\alpha '_2=(\\alpha _2+\\alpha_4\\xi_2) \\eta _2,$\n\n$\\alpha '_4=\\alpha _4 \\eta _2^2,$\n\n$\\beta '_1=\\frac{\\beta _1+(1+\\lambda)\\xi _2-\\alpha _4 \\xi _2^3}{\\eta _2}.$\n\\begin{itemize}\n\t\\item[1.] $\\alpha_4\\neq 0$. In this case we obtain\n\t$$ A_{2,3}(\\mathbf{c})=\\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1 & 0 & 0 & 1 \\\\\n\t\\beta _1& \\beta_2 & 1-\\alpha_1&0\n\t\\end{array}\\right)\\simeq \\left(\n\t\\begin{array}{cccc}\n\t\\alpha_1 & 0 & 0 & 1 \\\\\n\t-\\beta _1& \\beta_2 & 1-\\alpha_1&0\n\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1, \\beta_1, \\beta_2)\\in \\mathbf{F}^3.$$\n\t\n\t\\item[2.] $\\alpha _4=0.$\n\t\\begin{itemize}\n\t\t\\item[2-a)] If $\\alpha _2\\neq0$ one can make $\\alpha '_1=0$, $\\alpha'_2=1$ to get \\[ A_{3,3}(\\mathbf{c})=\\left(\n\t\t\\begin{array}{cccc}\n\t\t0 & 1 & 1 & 0 \\\\\n\t\t\\beta _1& \\beta_2 & 1&-1\n\t\t\\end{array}\\right)\\ \\mbox{with}\\ \\mathbf{c}=(\\beta_1,\\beta_2)\\in \\mathbf{F}^2.\\]\n\t\t\n\t\t\\item[2-b)] If $\\alpha _2=0$ then $\\alpha '_1=\\alpha _1,$ $\\alpha '_2=0,$ $\\alpha '_4=0,$ and $\\beta '_1=\\frac{\\beta _1+(1+\\lambda)\\xi _2}{\\eta _2}.$\n\t\t\\begin{itemize}\n\t\t\t\\item[2-b)-1.] Let $1+\\lambda\\neq 0$, that is $\\lambda\\neq -1$, one can make $\\beta'_1=0$ to get\n\t\t\t\\[A_{4,3}(\\mathbf{c})=\\left(\n\t\t\t\\begin{array}{cccc}\n\t\t\t\\alpha _1 & 0 & 0 & 0 \\\\\n\t\t\t0 & \\beta_2 & 1-\\alpha _1&0\n\t\t\t\\end{array}\\right), \\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1,\\beta_2)\\in \\mathbf{F}^2 \\ \\mbox{with}\\ \\beta_2\\neq -1-\\alpha_1.\\]\n\t\t\t\n\t\t\t\\item[2-b)-2.] Let $1+\\lambda=0.$\n\t\t\t\\begin{itemize}\n\t\t\t\t\\item If $\\beta_1\\neq 0$ one can make $\\beta'_1=1$ to get\n\t\t\t\t\\[A_{5,3}(\\mathbf{c})=\\left(\n\t\t\t\t\\begin{array}{cccc}\n\t\t\t\t\\alpha_1 & 0 & 0 & 0 \\\\\n\t\t\t\t1 & -1-\\alpha_1 & 1-\\alpha_1&0\n\t\t\t\t\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F}.\\]\n\t\t\t\t\\item If $\\beta_1=0$ one has $\\lambda=-1$ and therefore\n\t\t\t\t$A'=\\left(\n\t\t\t\t\\begin{array}{cccc}\n\t\t\t\t\\alpha _1 & 0 & 0 & 0 \\\\\n\t\t\t\t0 & -1-\\alpha _1 & 1-\\alpha _1&0\n\t\t\t\t\\end{array}\\right)$ which is $A_{4,3}(\\mathbf{c})$ with $\\beta_2=-1-\\alpha_1$.\n\t\t\t\\end{itemize}\n\t\t\\end{itemize}\n\t\\end{itemize}\n\\end{itemize}\n\n{\\bf The fourth subset case.} It is easy to see that in this case the result can be derived from the second and third subsets case. So we get\n\\[A_{6,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right)\\simeq \\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 1 \\\\\n-\\beta _1& 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right),\\ \\mbox{where}\\ \\mathbf{c}=(\\alpha_1,\\beta_1)\\in \\mathbf{F}^2\\]\\\\ (it is $A_{2,3}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\alpha_1+\\beta_2=0$),\n\\[ A_{7,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n0 & 1 & 1 & 0 \\\\\n\\beta_1& 1& 0&-1\n\\end{array}\\right), \\ \\mbox{where}\\ \\mathbf{c}=\\beta_1\\in \\mathbf{F}\\]\\\\ (it is $A_{3,3}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\beta_2=0$),\n\\[A_{8,3}(\\mathbf{c})=\\left(\n\\begin{array}{cccc}\n\\alpha_1 & 0 & 0 & 0 \\\\\n0 & 1-\\alpha_1 & -\\alpha_1&0\n\\end{array}\\right), \\mbox{where}\\ \\mathbf{c}=\\alpha_1\\in \\mathbf{F}\\]\\\\ (it is $A_{4,3}(\\mathbf{c})$, where the $2^{nd}$ and $3^{rd}$ columns are interchanged and $\\alpha_1+\\beta_2=0$).\n\n{\\bf The fifth subset case.}\n\nHere it is enough to consider $\\alpha_4=0$ case. Therefore we have\n\n$\\alpha' _1=\\frac{1}{\\Delta }\\left(-\\beta _1 \\eta _1 \\xi _1^2+\\alpha _1 \\eta _2 \\xi _1^2-\\alpha _1 \\eta _1 \\xi _1 \\xi _2-\\alpha _2 \\eta _2 \\xi _1 \\xi _2+\\alpha _2 \\eta _1 \\xi _2^2\\right),$\n\n$\\alpha' _2=\\frac{-1}{\\Delta }\\left(\\beta _1 \\eta _1^2 \\xi _1+\\alpha _1 \\eta _1 \\eta _2 \\xi _1-\\alpha _2 \\eta _2^2 \\xi _1-\\alpha _1 \\eta _1^2 \\xi _2+\\alpha _2 \\eta _1 \\eta _2 \\xi _2\\right),$\n\n$\\alpha' _4=\\frac{-1}{\\Delta }\\beta _1 \\eta _1^3,$\n\n$\\beta' _1=\\frac{1}{\\Delta }\\beta _1 \\xi _1^3.$\n\nMaking $\\eta_1=0$ results in $\\alpha'_1=\\alpha_1 \\xi_1-\\alpha _2\\xi_2,$ $\\alpha'_4=0$, $\\alpha'_2=\\alpha _2\\eta_2$, $\\beta'_1=\\frac{\\xi^2_1}{\\eta_2}\\beta _1.$\n\\begin{itemize}\n\t\\item[a)] If $\\alpha_2\\neq 0$ one can make $\\alpha'_1= 0$ and $\\alpha'_2=1$.\n\t\\begin{itemize}\n\t\t\\item[a)-1.] If $\\beta_1\\neq 0$ one can reduce $\\beta'_1=1$ to get\n\t\t$A_{9,3}=\\left(\n\t\t\\begin{array}{cccc}\n\t\t0 & 1 & 1 & 0 \\\\\n\t\t1 &0&0 &-1\\end{array}\n\t\t\\right).$\n\t\t\n\t\t\\item[a)-2.] If $\\beta_1= 0$ then $\\beta'_1=0$ and one gets\n\t\t$A_{10,3}=\\left(\n\t\t\\begin{array}{cccc}\n\t\t0 &1 & 1 & 0 \\\\\n\t\t0 &0& 0 &-1\n\t\t\\end{array}\n\t\t\\right).$\n\t\\end{itemize}\n\t\\item[b)] If $\\alpha_2= 0$ then $\\alpha'_2= 0$.\n\t\\begin{itemize}\n\t\t\\item[b)-1.] If $\\beta_1\\neq 0$ we reduce $\\beta'_1=1$ and $\\alpha'_1$ may take one or zero depending on $\\alpha_1$. Therefore one has\n\t\t\\[A_{11,3}=\\left(\n\t\t\\begin{array}{cccc}\n\t\t1 &0 & 0 & 0 \\\\\n\t\t1 &-1&-1 &0\n\t\t\\end{array}\n\t\t\\right)\\ \\mbox{or}\\ A_{12,3}=\\left(\n\t\t\\begin{array}{cccc}\n\t\t0 &0 &0 & 0 \\\\\n\t\t1 &0& 0 &0\n\t\t\\end{array}\n\t\t\\right).\\]\n\t\t\n\t\t\\item[b)-2.] If $\\beta_1=0$ then $\\beta'_1=0$, $\\alpha'_4=0$ and one can make , $\\alpha'_1=1$ to get\n\t\t\\[A'=\\left(\n\t\t\\begin{array}{cccc}\n\t\t1 &0& 0 & 0 \\\\\n\t\t0 &-1& -1 &0\n\t\t\\end{array}\n\t\t\\right)\\ \\mbox{isomorphic to}\\ A_{10,3}.\\]\n\t\\end{itemize}\n\\end{itemize}\n\n{\\bf Remark 1.} From the proof of Theorems it is clear that the same results remain be true if one assumes the existence of a root in $\\mathbf{F}$ of every second and third order polynomial over $\\mathbf{F}$ (instead of assumption that $\\mathbf{F}$ to be algebraically closed). The results can be used for getting complete classification of different classes, for example, associative, Jordan, alternative and etc., 2-dimensional algebras. \n\n{\\bf Remark 2.} In \\cite{A} the class $A_{3,2}(\\mathbf{c})$ should be understood as it is in this paper.\n \n{\\bf Acknowledgment.}\n The first author thanks M.A.A. Ahmed for the help in computer computations, the second author's research is supported by FRGS14-153-0394, MOHE and the third author acknowledges MOHE for a support by grant 01-02-14-1591FR.\n\n\n\\vspace{0,2cm}\n\n\\bibliographystyle{aipproc}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section*{Introduction}\n\nIn this paper, we prove the conjecture by Gyenge, N\\'emethi and\nSzendr\\H{o}i \\cite{MR3671513,MR3799153} giving a formula of the\ngenerating function of Euler numbers of Hilbert schemes of points\n$\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$ on a simple singularity\n${\\mathbb C}^2\/\\Gamma$, where $\\Gamma$ is a finite subgroup of $\\operatorname{\\rm SL}(2)$. When\n$\\Gamma$ is of type $A$, Euler numbers were computed by Dijkgraaf and\nSu\\l kowski \\cite{MR2391107}, and Toda \\cite{MR3394393}. The formula in\n\\cite{MR3671513,MR3799153} is given in a different form, and makes\nsense for arbitrary $\\Gamma$. \nThe formula was proved for type $D$, as well as type $A$, in\n\\cite{MR3799153}.\nOur proof of the formula in type $E$ is new.\n\nThe formula \\cite{MR3671513,MR3799153} is written as a specialization\nof the character for the `Fock space', which is the tensor product of\nthe basic representation of the affine Lie algebra ${\\mathfrak{g}}_{\\mathrm{aff}}$\ncorresponding to $\\Gamma$, and the usual Fock space representation of\nthe Heisenberg algebra. The Fock space was realized in a geometric way\nas the direct sum of homology groups of various connected components\nof $\\operatorname{Hilb}^n({\\mathbb C}^2)^\\Gamma$ for all $n$. See\n\\cite{MR1989196}. The character is well-known (see e.g.,\n\\cite[\\S12]{Kac})\n\\begin{equation*}\n Z =\n \\prod_{m=1}^\\infty\n (1 - e^{-m\\delta})^{-(n+1)}\n \\sum_{\\vec{m}\\in{\\mathbb Z}^n} \n e^{-\\frac{(\\vec{m}, \\mathbf C\\vec{m})}2\\delta}\n \\prod_{i=1}^n e^{-m_i\\alpha_i},\n\\end{equation*}\nwhere $\\alpha_i$ is the $i$th simple root of ${\\mathfrak{g}}_{\\mathrm{fin}}$, the\nunderlying finite dimensional complex simple Lie algebra, $\\delta$ is\nthe positive primitive imaginary root, $n$ is the rank of\n${\\mathfrak{g}}_{\\mathrm{fin}}$, and $\\mathbf C$ is the finite Cartan matrix for\n${\\mathfrak{g}}_{\\mathrm{fin}}$. Here we omit the factor $e^{\\Lambda_0}$ for the $0$th\nfundamental weight $\\Lambda_0$. The well-known formula for the basic\nrepresentation gives $(1 - e^{-m\\delta})^{-n}$ instead of the power\n$-(n+1)$. The additional $-1$ comes from the Heisenberg algebra.\n\\begin{NB}\n Suppose the dimension vector $\\bv$ for ${\\mathfrak{M}}(\\bv,\\Lambda_0)$ is\n $m\\delta + \\bv_{\\mathrm{fin}}$ with\n $\\bv_{\\mathrm{fin}} = \\sum_{i\\neq 0} m_i \\alpha_i$. By the Frenkel-Kac\n construction, we have the `bottom' case with\n $\\mathop{\\text{\\rm dim}}\\nolimits {\\mathfrak{M}}(\\bv,\\Lambda_0) = 0$, and raise its vacuum by the\n Heisenberg algebra. For the bottom, we have\n \\begin{equation*}\n 0 = \\mathop{\\text{\\rm dim}}\\nolimits {\\mathfrak{M}}(m\\delta+\\bv_{\\mathrm{fin}},\\Lambda_0)\n = 2m - (m\\delta+\\bv_{\\mathrm{fin}}, \\mathbf C(m\\delta+\\bv_{\\mathrm{fin}}))\n = 2m - (\\bv_{\\mathrm{fin}}, \\mathbf C \\bv_{\\mathrm{fin}})\n \\end{equation*}\n as $\\mathbf C\\delta = 0$. Thus\n $m = \\frac12 (\\bv_{\\mathrm{fin}}, \\mathbf C \\bv_{\\mathrm{fin}})$.\n\\end{NB}%\n\nThe conjectural formula in \\cite{MR3671513}, which we will prove, is\nthe following:\n\\begin{Theorem*}\\label{thm:1}\n The generating function of Euler numbers of\n $\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$\n \\begin{equation*}\n \\sum_{n=0}^\\infty e^{-n\\delta} \\chi(\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma))\n \\end{equation*}\n is obtained from $Z$ by the substitution\n \\begin{equation*}\n e^{-\\alpha_i} = \\exp\\left( \\frac{2\\pi i}{h^\\vee+1}\\right)\n \\quad (i\\neq 0),\n \\end{equation*}\n where $h^\\vee$ is the dual Coxeter number.\n\\end{Theorem*}\n\nLet us sketch the strategy of our proof.\nIn order to compute Euler numbers of\n$\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$, we use a recent result of Craw,\nGammelgaard, Gyenge and Szendr\\H{o}i \\cite{2019arXiv191013420C}:\n$\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$ with the reduced scheme\nstructure is isomorphic to a quiver variety\n${\\mathfrak{M}}_{\\zeta^\\bullet}(n\\delta,\\Lambda_0)$, where $\\zeta^\\bullet$ is a\nstability parameter\n\\begin{equation*}\n \\zeta_i^\\bullet = 0 \\quad\\text{for all $i\\neq 0$},\\quad\n \\zeta_0^\\bullet < 0.\n\\end{equation*}\nThis stability parameter is different from the standard one $\\zeta$\n\\begin{equation*}\n \\zeta_i < 0\\quad\\text{for all $i\\in I$},\n \n \n\\end{equation*}\nwhich gives connected components ${\\mathfrak{M}}_{\\zeta}(\\bv, \\Lambda_0)$ of\n$\\operatorname{Hilb}^{n}({\\mathbb C}^2)^\\Gamma$. Here $\\bv$ is an isomorphism\nclass of ${\\mathbb C}[x,y]\/I$ as a $\\Gamma$-module for an ideal\n$I \\subset {\\mathbb C}[x,y]$ of the coordinate ring ${\\mathbb C}[x,y]$ of ${\\mathbb C}^2$,\ninvariant under the $\\Gamma$-action.\n\nAs $\\zeta^\\bullet$ lies in walls given by roots in the space of\nstability parameters, ${\\mathfrak{M}}_{\\zeta^\\bullet}(n\\delta,\\Lambda_0)$ is\nsingular in general. Nevertheless it has a representation theoretic\nmeaning, as shown in \\cite{Na-branching}: the quiver variety\n${\\mathfrak{M}}_{\\zeta^\\bullet}(n\\delta,\\Lambda_0)$ is responsible for the\nrestriction of ${\\mathfrak{g}}_{\\mathrm{aff}}$ to ${\\mathfrak{g}}_{\\mathrm{fin}}$.\nThe essential geometric ingredient for this relation is a projective\nmorphism\n\\begin{equation*}\n \\pi_{\\zeta^\\bullet,\\zeta}\\colon\n {\\mathfrak{M}}_{\\zeta}(\\bv,\\Lambda_0)\n \\to {\\mathfrak{M}}_{\\zeta^\\bullet}(n\\delta,\\Lambda_0)\n = \\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma),\n \\qquad n = \\bv_0.\n\\end{equation*}\nBy a local description of singularities of quiver varieties\n(\\cite[\\S2.7]{Na-branching}, which went back to\n\\cite[\\S6]{Na-quiver}), fibers of $\\pi_{\\zeta^\\bullet,\\zeta}$ are\nlagrangian subvarieties in quiver varieties ${\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)$\nassociated with the finite $ADE$ quiver for ${\\mathfrak{g}}_{\\mathrm{fin}}$, the finite\ndimensional complex simple Lie algebra underlying ${\\mathfrak{g}}_{\\mathrm{aff}}$.\nAs a simple application of this result, we observe that\n\\begin{itemize}\n\\item Euler numbers of $\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$ can be\n written in terms of Euler numbers of ${\\mathfrak{M}}_\\zeta(\\bv,\\Lambda_0)$ and\n quiver varieties ${\\mathfrak{M}}_\\zeta(\\bv^s,\\bw^s)$ of the finite $ADE$ type.\n\\end{itemize}\nIn fact, we get \\emph{more} data than we need: we have a\nstratification of $\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$ so that\n$\\pi_{\\zeta^\\bullet,\\zeta}$ is a fiber bundle over each stratum. We\ncan compute Euler numbers of all strata from Euler numbers of\n${\\mathfrak{M}}_\\zeta(\\bv,\\Lambda_0)$ and quiver varieties\n${\\mathfrak{M}}_\\zeta(\\bv^s,\\bw^s)$ of the finite $ADE$ type. See \\eqref{eq:5}\nfor the precise formula.\n\nThere are several algorithms\\footnote{We make a distinction between an\n algorithm for a computation and an explicit computation as in\n \\cite{Na:E8}. Namely a computation results finitely many $\\pm$,\n $\\times$, integers and variables. We do not require that the final\n expression is readable by a human. The expression like\n $\\sum_{i=1}^{2^{(2^{100})}} a_i$ with explicit $a_i$ is an\n algorithm.} to compute Euler numbers of quiver varieties. For\nexample, \\cite{MR2144973} and \\cite{MR2651380}. They are complicated,\nand hard to use in practice.\nOn the other hand, the conjectural formula in \\cite{MR3671513} above\ndoes \\emph{not} contain such a complicated algorithm. It is just given\nby a simple substitution.\nIt means that we should have a drastic simplification if we take a\n\\emph{linear combination} of complicated Euler numbers of quiver\nvarieties of finite $ADE$ type. We do not need to compute Euler\nnumbers of individual strata, as we only need their sum.\n\nThis simplification has a representation theoretic origin. Let us\nexplain it.\n\nThe above specialized character is called the \\emph{quantum dimension}:\n\\begin{equation*}\n \\mathop{\\text{\\rm dim}}\\nolimits_q V = \\left.\\operatorname{ch} V\\right|_{e^{-\\alpha_i}=\\exp(\\frac{2\\pi i}{h^\\vee+1})}.\n\\end{equation*}\nHere $V$ is a finite dimensional representation $V$ of ${\\mathfrak{g}}_{\\mathrm{fin}}$.\nIt was introduced by Andersen~\\cite[Def.~3.1]{MR1182414} (see also\nParshall-Wang~\\cite{MR1247547}). We choose a specific root of unity\n$\\exp(\\frac{2\\pi i}{2(h^\\vee+1)})$, where these papers study more\ngeneral roots of unity $\\zeta$.\n\nTherefore the conjectural formula in \\cite{MR3671513} states that the\ngenerating function of Euler numbers of\n$\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$ is given by the quantum\ndimension of the Fock space at $\\exp(\\frac{2\\pi i}{2(h^\\vee+1)})$,\nrestricted from ${\\mathfrak{g}}_{\\mathrm{aff}}$ to ${\\mathfrak{g}}_{\\mathrm{fin}}$, just keeping track\n$e^{-\\delta}$.\n\nA direct sum of homology groups (which are isomorphic to complexified\n$K$-group) of Lagrangian subvarieties of quiver varieties of the\nfinite $ADE$ type carries a structure of a finite dimensional\nrepresentation of the quantum loop algebra\n$\\mathbf U_{q}(\\mathbf L{\\mathfrak{g}}_{\\mathrm{fin}})$. See \\cite{Na-qaff}. It is called a\n\\emph{standard module} of $\\mathbf U_{q}(\\mathbf L{\\mathfrak{g}}_{\\mathrm{fin}})$.\n\nIt turns out that the above simplification is a consequence of the\nfollowing representation theoretic result of an independent interest:\n\\begin{Theorem*}\\label{thm:2}\n The quantum dimention of arbitrary standard module of\n $\\mathbf U_q(\\mathbf L{\\mathfrak{g}}_{\\mathrm{fin}})$ of type $ADE$ at\n $\\zeta=\\exp(\\frac{2\\pi i}{2(h^\\vee+1)})$ is equal to $1$.\n\\end{Theorem*}\n\nSince standard modules are tensor products of $l$-fundamental modules\n\\cite{VV-std}, it is enough to prove this result for $l$-fundamental\nmodules.\nThe author is told by Naoi that this result is a special case of more\ngeneral conjecture posed by Kuniba \\cite[Conj.~2 (A.6a)]{MR1202213}\n(see also Kuniba-Nakanishi-Suzuki \\cite[Conjecture~14.2]{MR2773889})\nformulated for Kirillov-Reshetikhin modules.\n(Recall $l$-fundamental modules are the simplest examples of\nKirillov-Reshetikhin modules.)\nIt is not difficult to check the general conjecture for type $A$. Type\n$D$ case was shown in \\cite{MR3043892}, while type $E_6$ case was\nshown in \\cite{MR3282650}.\n\nAlthough we only need the \\emph{simplest special} case of more general\nconjecture, we could not find a proof of the relevant result for type\n$E_7$, $E_8$ in the literature. Therefore we give its proof.\nFortunately the necessary explicit computation for type $E$ from the\nalgorithm in \\cite{MR2144973} was already done in \\cite{Na:E8} by\nusing a supercomputer. Alternatively we could quote the computation by\nKleber \\cite{MR1436775}, which assumed fermionic formula conjectured\nat that time. The fermionic formula was proved later by\nDi~Francesco and Kedem~\\cite{MR2428305}.\n\\begin{NB}\n based on many earlier works including \\cite{MR1993360}, as well as\n results from cluster algebras.\n\\end{NB}%\n\nWe also give a proof of the known cases $A$, $D$, $E_6$ for the\ncompleteness. We encounter a new feature in $E_7$, $E_8$, which did\nnot arise in other cases. Hence our check in the simplest case is yet\nnontrivial.\n\nBecause it is about Kirillov-Reshetikhin modules, it suggests that a\ngeneral conjecture should be studied in the framework of cluster\nalgebras, as in \\cite{MR2428305}. Note also that Euler numbers of\n(graded) quiver varieties are understood in the context of cluster\nalgebras in a recent work of Bittmann \\cite{2019arXiv191113110B}. Thus\nthe suggestion is compatible with the approach in this paper, though\ncluster algebras play no role in this paper.\n\nThe paper is organized as follows. In \\cref{sec:Euler} we deduce\n\\cref{thm:1} from \\cref{thm:2} after recalling results from\n\\cite{Na-branching}. In \\cref{sec:2} we prove \\cref{thm:2}. It is\nproved by the case by case analysis. In \\cref{sec:rationally} we\ndiscuss an additional topic, the rationally smoothness of\n$\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$. It is rationally smooth if\n$n=2$, but not so in general.\n\n\\subsection*{Acknowledgments}\n\nThe author thanks\nA.~Gyenge\nand\nB.~Szendr\\H{o}i\nfor discussion on \\cite{MR3671513,MR3799153,2019arXiv191013420C}, and\nA.~Craw\non \\cite{2019arXiv191013420C} and \\cite{MR1202213,MR2773889}.\nHe also thanks\nA.~Kuniba,\nS.~Naito,\nT.~Nakanishi\nand\nK.~Naoi\nfor discussion on the conjecture in \\cite{MR1202213,MR2773889},\nD.~Muthiah\nfor discussion on rational smoothness,\nand\nR.~Yamagishi\nfor discussion on \\cite{2017arXiv170905886Y}.\n\nThe research is supported in part by the World Premier International\nResearch Center Initiative (WPI Initiative), MEXT, Japan, and by JSPS\nGrants Numbers 16H06335, 19K21828.\n\n\\section{Euler numbers of quiver varieties}\\label{sec:Euler}\n\n\\subsection{Quiver varieties}\n\nLet $\\Gamma$ be a nontrivial finite subgroup of $\\operatorname{\\rm SL}(2)$.\nWe define the affine Dynkin diagram via the McKay correspondence (see\ne.g., \\cite[Ch.~8]{MR3526103} for detail): let\n$\\{ \\rho_i\\}_{i\\in I}$ be the set of isomorphism classes of\nirreducible representations of $\\Gamma$ with the trivial\nrepresentation $\\rho_0$. We identify $\\rho_i$ with a vertex of a\ngraph. We draw $a_{ij}$ edges between $\\rho_i$ and $\\rho_j$ where\n$a_{ij} = \\mathop{\\text{\\rm dim}}\\nolimits \\operatorname{Hom}_\\Gamma(\\rho_i, \\rho\\otimes\\rho_j) =\n\\mathop{\\text{\\rm dim}}\\nolimits\\operatorname{Hom}_\\Gamma(\\rho_j,\\rho\\otimes\\rho_i)$, where $\\rho$ is the\n$2$-dimensional representation of $\\Gamma$ given by the inclusion\n$\\Gamma\\subset\\operatorname{\\rm SL}(2)$. Then the graph is an affine Dynkin diagram of\ntype $ADE$. Let ${\\mathfrak{g}}_{\\mathrm{aff}}$ denote the corresponding affine Lie\nalgebra, and ${\\mathfrak{g}}_{\\mathrm{fin}}$ the underlying finite dimensional complex\nsimple Lie algebra corresponding to the Dynkin diagram obtained from\nthe affine one by removing $\\rho_0$. Let $n$ be the rank of\n${\\mathfrak{g}}_{\\mathrm{fin}}$, which is the number of vertices in $I$ minus $1$.\n\nWe use the convention of the root system for ${\\mathfrak{g}}_{\\mathrm{aff}}$ as in\n\\cite[Ch.~6 and \\S12.4]{Kac}.\nLet $\\alpha_i$ be the $i$th simple root of ${\\mathfrak{g}}_{\\mathrm{aff}}$ corresponding\nto $\\rho_i$, $\\delta$ be the primitive positive imaginary root of\n${\\mathfrak{g}}_{\\mathrm{aff}}$. We have $\\delta = \\sum a_i \\alpha_i$, and $a_i$ is equal\nto the dimension of $\\rho_i$.\nLet $\\alpha_i^\\vee$ be the $i$th simple coroot. We take the scaling\nelement $d$ satisfying\n\\begin{equation*}\n \\langle\\alpha_i, d\\rangle = \\delta_{i0}.\n\\end{equation*}\nThen $\\langle \\alpha_i^\\vee, d\\rangle_{i\\in I}$ is a base of\n$\\mathfrak h_{\\mathrm{aff}}$, the Cartan subalgebra of ${\\mathfrak{g}}_{\\mathrm{aff}}$.\nWe define the fundamental weights $\\Lambda_i$ ($i\\in I$) by\n\\begin{equation*}\n \\langle \\Lambda_i, \\alpha_j^\\vee\\rangle = \\delta_{ij},\\quad\n \\langle \\Lambda_i, d\\rangle = 0.\n\\end{equation*}\nThen $\\langle \\alpha_i, \\Lambda_0\\rangle_{i\\in I}$ forms a base of\n$\\mathfrak h_{\\mathrm{aff}}^*$.\n\nWe choose an orientation of edges in the affine Dynkin diagram and\nconsider the corresponding affine quiver $Q = (I,\\Omega)$.\nWe take dimension vectors $\\bw = (\\bw_i)$, and\n$\\bv = (\\bv_i)\\in {\\mathbb Z}_{\\ge 0}^{I}$, and\nconsider quiver varieties ${\\mathfrak{M}}_\\zeta(\\bv,\\bw)$,\n${\\mathfrak{M}}_{\\zeta^\\bullet}(\\bv,\\bw)$, where $\\zeta$, $\\zeta^\\bullet$ are\nstability parameters such that\n\\begin{equation*}\n \\zeta_i < 0\\quad\\text{for all $i\\in I$},\\qquad\n \\zeta_i^\\bullet = 0 \\quad\\text{for all $i\\neq 0$},\\quad\n \\zeta_0^\\bullet < 0.\n\\end{equation*}\nSee \\cite[\\S2]{Na-branching} for the definition of quiver varieties\nfor these stability conditions.\nSince $\\zeta^\\bullet$ lives in the boundary of a chamber containing\n$\\zeta$, we have a projective morphism\n\\begin{equation*}\n \\pi_{\\zeta^\\bullet,\\zeta}\\colon\n {\\mathfrak{M}}_\\zeta(\\bv,\\bw)\\to {\\mathfrak{M}}_{\\zeta^\\bullet}(\\bv,\\bw).\n\\end{equation*}\nSee \\cite[\\S2]{Na-branching}.\nIn fact, this is an example studied in \\cite[\\S2.8]{Na-branching}\nassociated with a division $I = I^0\\sqcup I^+$ with\n$I^0 = I\\setminus \\{0\\}$, $I^+ = \\{0\\}$.\n\nWe identify $\\bv$, $\\bw$ with weights of ${\\mathfrak{g}}_{\\mathrm{aff}}$ by\n\\begin{equation*}\n \\bv = \\sum \\bv_i \\alpha_i, \\quad \\bw = \\sum \\bw_i \\Lambda_i.\n\\end{equation*}\nIt is convenient to use a different convention for the dimension\nvector $\\bv$:\n\\begin{equation}\\label{eq:3}\n \\bv = m\\delta + \\sum_{i\\in I^0} m_i \\alpha_i, \\qquad \\text{i.e., }\n m = \\bv_0, \\quad m_i = \\bv_i - \\bv_0 a_i.\n\\end{equation}\n\nLet us take $\\bw = \\Lambda_0$.\nIt is well-known that ${\\mathfrak{M}}_\\zeta(\\bv,\\Lambda_0)$ is the $\\Gamma$-fixed\npoint component of Hilbert schemes $I$ of points in the affine plane\n${\\mathbb C}^2$ such that ${\\mathbb C}[x,y]\/I$ is isomorphic to\n$\\bigoplus \\rho_i^{\\oplus \\bv_i}$, as a $\\Gamma$-module.\nEuler number for ${\\mathfrak{M}}_\\zeta(\\bv,\\Lambda_0)$ is also known. It was\ngiven in \\cite{MR1989196}. Since it was stated without an explanation,\nlet us explain how it is derived. We change the stability condition\n$\\zeta$ to see that ${\\mathfrak{M}}_\\zeta(\\bv,\\bw)$ is diffeomorphic to a moduli\nspace of framed rank $1$ torsion free sheaves on the minimal\nresolution of ${\\mathbb C}^2\/\\Gamma$. Then rank $1$ torsion free sheaves are\nideal sheaves twisted by line bundles, hence Euler numbers are given\nby G\\\"ottsche formula for Hilbert schemes of points \\cite{MR1032930}.\nMoreover this latter picture gives the Frenkel-Kac construction of the\nbasic representation of ${\\mathfrak{g}}_{\\mathrm{aff}}$ (see e.g., \\cite[\\S14.8]{Kac}),\nhence\nwe get the formula of $Z$ in Introduction. We have used the convention\nin \\eqref{eq:3}.\n\nWe can define a structure of a representation of the affine Lie\nalgebra ${\\mathfrak{g}}_{\\mathrm{aff}}$ on the direct sum of homology groups\n$\\bigoplus_\\bv H_*({\\mathfrak{M}}_\\zeta(\\bv,\\Lambda_0))$. See \\cite{MR1989196}\nand references therein. We can also construct a structure of a\nrepresentation of the Heisenberg algebra commuting with ${\\mathfrak{g}}_{\\mathrm{aff}}$\nby \\cite{MR1441880,Lecture}.\n\nOn the other hand, Craw, Gammelgaard, Gyenge and Szendr\\H{o}i\n\\cite{2019arXiv191013420C} recently proved that the Hilbert scheme\n$\\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$ of $n$ points in ${\\mathbb C}^2\/\\Gamma$\nwith the reduced scheme structure is isomorphic to\n${\\mathfrak{M}}_{\\zeta^\\bullet}(\\bv,\\Lambda_0)$ with $\\bv = n\\delta$.\n\n\\subsection{Stratification}\\label{subsec:strat}\n\nAs is explained in \\cite[\\S2]{Na-branching},\n${\\mathfrak{M}}_{\\zeta^\\bullet}(\\bv,\\bw)$ parametrizes $S$-equivalence classes of\n$\\zeta^\\bullet$-semistable framed representations of the preprojective\nalgebra. Therefore its points are represented by direct sum of\n$\\zeta^\\bullet$-stable representations. Under the above choice of\n$\\zeta^\\bullet$, we have one distinguished summand, giving a point in\n${\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\\bv',\\bw)$ for some $\\bv'\\le\\bv$\n(component-wise) such that $\\bv'_0 = \\bv_0$, and other summands are\nsimple representations $S_i$ with $i\\in I^0$. See\n\\cite[\\S2.6]{Na-branching}.\nSince multiplicities of $S_i$ can be read off from the difference\n$\\bv - \\bv'$, we can regard\n${\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\\bv',\\bw)$ as a subset of\n${\\mathfrak{M}}_{\\zeta^\\bullet}(\\bv,\\bw)$.\nBy \\cite[Prop.~2.30]{Na-branching}\n${\\mathfrak{M}}^{\\mathrm{s}}_{\\zeta^\\bullet}(\\bv',\\bw)\\neq\\emptyset$ if and only\nif $\\bw - \\bv'$ is an $I^0$-dominant weight appearing in the basic\nrepresentation $V(\\Lambda_0)$ of the affine Lie algebra ${\\mathfrak{g}}$\nassociated with $\\Gamma$.\nHere $I^0$-dominant means that\n$\\langle \\bw - \\bv', \\alpha_i^\\vee\\rangle \\ge 0$ for $i\\in I^0$.\nWe thus have the stratification\n\\begin{equation*}\n {\\mathfrak{M}}_{\\zeta^\\bullet}(\\bv,\\bw)\n = \\bigsqcup_{\\text{$\\bv'$ as above}} {\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\\bv',\\bw). \n\\end{equation*}\n\nMoreover the transversal slice to the stratum\n${\\mathfrak{M}}^{\\mathrm{s}}_{\\zeta^\\bullet}(\\bv',\\bw)$ is locally isomorphic to\na quiver variety ${\\mathfrak{M}}_0(\\bv^s,\\bw^s)$ around $0$, associated with the\nfinite $ADE$ quiver $Q\\setminus \\{\\rho_0\\}$ such that dimension\nvectors are given by\n\\begin{equation*}\n \\bv^s = \\bv - \\bv', \\qquad\n \\bw^s_i = \\langle \\bw - \\bv', \\alpha_i^\\vee\\rangle \\quad i\\in I^0.\n\\end{equation*}\nNote that $\\bv-\\bv'$ has no $0$th component as $\\bv'_0 = \\bv_0$.\nNote also that $\\bw^s_i \\ge 0$, as $\\bw - \\bv'$ is $I^0$-dominant.\nIt is also known that the inverse image of the slice in\n${\\mathfrak{M}}_\\zeta(\\bv,\\bw)$ under $\\pi_{\\zeta^\\bullet,\\zeta}$ is locally\nisomorphic to ${\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)$ around\n$\\mathfrak L(\\bv^s,\\bw^s)$, the inverse image of the origin $0$ of\n${\\mathfrak{M}}_0(\\bv^s,\\bw^s)$ under the projective morphism\n${\\mathfrak{M}}_\\zeta(\\bv^s,\\bw^s)\\to {\\mathfrak{M}}_0(\\bv^s,\\bw^s)$.\nSee \\cite[\\S2.7]{Na-branching} and the references therein for these\nclaims on transversal slices. There is also an algebraic approach in\n\\cite{CB:normal}.\n\nFor $\\bw = \\Lambda_0$, we have\n\\begin{equation*}\n \\bw^s_i = - \\sum_{j\\in I^0} m'_j \\langle \\alpha_j, \\alpha_i^\\vee\\rangle,\n \\qquad \\text{where $\n \\bv' = m\\delta + \\sum_{j\\in I^0} m'_j \\alpha_j$\n in the convention \\eqref{eq:3}}.\n\\end{equation*}\nNote that $\\bv = m\\delta + \\sum m_i \\alpha_i$ and $\\bv'$ share the\nsame $m$ for the coefficient of $\\delta$, as $\\bv_0 = \\bv'_0$.\nIn particular,\n\\begin{equation}\\label{eq:6}\n \\sum_{i\\in I^0} \\bw^s_i \\Lambda_i - \\bv^s_i \\alpha_i\n \\begin{NB}\n = - \\sum_{i\\in I^0} \\sum_{j\\in I^0} m'_j \\langle\\alpha_j,\n \\alpha_i^\\vee \\rangle \\Lambda_i - (m_i - m'_i)\\alpha_i\n \\end{NB}\n = - \\sum_{i\\in I^0} m_i \\alpha_i.\n\\end{equation}\nWe also note\n\\begin{equation}\\label{eq:7}\n m'_i \\le 0 \\qquad (i\\in I^0),\n\\end{equation}\nas $(m'_i)_{i\\in I^0} = - (\\bw^s)_{i\\in I^0} \\mathbf C^{-1}$, and the\ninverse of the Cartan matrix $\\mathbf C$ has positive entries. This\nwas proved in \\cite[Prop.~A.1]{2019arXiv191013420C} by a different\nmethod.\n\n\\subsection{Examples}\n\n\\begin{Example}\n Let\n ${\\mathfrak{M}}_{\\zeta^\\bullet}(n\\delta, \\Lambda_0)\\cong\n \\operatorname{Hilb}^n({\\mathbb C}^2\/\\Gamma)$.\n Consider a stratum containing $(n-1)$-distinct points in\n ${\\mathbb C}^2\\setminus\\{0\\}\/\\Gamma$ together with trivial\n representations. We have\n $\\bv' = (n-1)\\delta + \\alpha_0 = n\\delta - \\sum_{i\\in I^0} a_i\n \\alpha_i$, and hence\n \\begin{equation*}\n \\bv^s = \\sum_{i\\in I^0} a_i \\alpha_i, \\quad\n \\bw^s_i = \\sum_{j\\in I^0} C_{ij} a_j.\n\\end{equation*}\nNote that $\\bw^s$ has entries $1$ at vertices in the finite quiver\n$I^0$ which are connected to the $0$-vertex in the affine quiver. In\nparticular, ${\\mathfrak{M}}_0(\\bv^s,\\bw^s)$ is ${\\mathbb C}^2\/\\Gamma$,\n${\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)$ is its minimal resolution, the very\nfirst example of a quiver variety considered by Kronheimer \\cite{Kr},\nbefore the definition of quiver varieties was introduced.\n\nThis is obvious, as the transversal slice is\n${\\mathfrak{M}}_{\\zeta^\\bullet}(\\delta,\\Lambda_0) \\cong {\\mathbb C}^2\/\\Gamma$, as we can\nignore $(n-1)$-distinct points.\n\\end{Example}\n\n\\begin{Example}[Yamagishi \\cite{2017arXiv170905886Y}]\\label{ex:Yam}\n Consider\n ${\\mathfrak{M}}_{\\zeta^\\bullet}(2\\delta, \\Lambda_0) =\n \\operatorname{Hilb}^2({\\mathbb C}^2\/\\Gamma)$. The formal neighborhoods of\n fibers of $\\pi_{\\zeta^\\bullet,\\zeta}$ over $0$-dimensional strata in\n ${\\mathfrak{M}}_{\\zeta^\\bullet}(2\\delta,\\Lambda_0)$ were determined by Yamagishi\n \\cite{2017arXiv170905886Y}. He identified the formal neighborhood\n with that of the intersection of the nilpotent cone for the complex\n simple Lie algebra ${\\mathfrak{g}}_{\\mathrm{fin}}$ and the Slodowy slice to a\n `sub-subregular' orbit. There is only one $0$-dimensional stratum\n except type $A_1$, $A_2$, $D_n$, while there are none for $A_1$,\n $A_2$, two for $D_n$ ($n > 4$) and three for $D_4$.\n We can give corresponding vector $\\bv'$ as follows:\n \\begin{itemize}\n \\item If $Q$ is of type $A_1^{(1)}$ or $A_2^{(1)}$, there is no such\n $\\bv'$.\n \\item If $Q$ is of type $A_{n}^{(1)}$ with $n > 2$,\n $\\bv' = 2\\alpha_0 + \\alpha_1 + \\alpha_{n}$.\n \\item If $Q$ is of type $D_n^{(1)}$ with $n > 4$,\n $\\bv' = \\dynkin[edge\n length=.6cm,labels={2,0,2,2,2,2,1,1}]{D}[1]{***.****} = 2\\alpha_0\n + 2\\alpha_2 + \\dots + 2\\alpha_{n-2} + \\alpha_{n-1} + \\alpha_n$ or\n $\\bv' = \\dynkin[edge\n length=.6cm,labels={2,1,2,1,0,0,0,0,0}]{D}[1]{****.****} =\n 2\\alpha_0 + \\alpha_1 + 2\\alpha_2 + \\alpha_3$.\n \\item If $Q$ is of type $D_4^{(1)}$, we have three possibilities: in\n addition to the second example in $D_n^{(1)}$ above, we have\n $(\\alpha_1,\\alpha_3) \\mapsto (\\alpha_1,\\alpha_4)$,\n $(\\alpha_3,\\alpha_4)$, as its cyclic permutation.\n\n \\item If $Q$ is of type $E_8^{(1)}$, we have\n $\\bv' = \\dynkin[edge\n length=.6cm,labels={2,0,1,1,2,2,2,2,2}]{E}[1]{8}$. The cases\n $E_6^{(1)}$, $E_7^{(1)}$ are similar.\n \\end{itemize}\n Here $0$th vertex is $\\circ$, and other vertices are $\\bullet$.\n\n Transversal slices are\n \\begin{itemize}\n \\item If $Q$ is of type $A_{n}^{(1)}$ with $n>2$,\n $\\bv^s = \\dynkin[edge\n length=.6cm,labels={1,2,2,2,2,1}]{A}{***.***}$,\n $\\bw^s = \\dynkin[edge\n length=.6cm,labels={0,1,0,0,1,0}]{A}{***.***}$.\n\n \\item If $Q$ is of type $D_n^{(1)}$,\n $\\bv^s = \\dynkin[edge\n length=.6cm,labels={2,2,2,2,1,1}]{D}{***.***}$,\n $\\bw^s = \\dynkin[edge\n length=.6cm,labels={2,0,0,0,0,0}]{D}{***.***}$ in the first case,\n and\n $\\bv^s = \\dynkin[edge\n length=.6cm,labels={1,2,3,4,4,4,4,2,2}]{D}{*****.****}$,\n $\\bw^s = \\dynkin[edge\n length=.6cm,labels={0,0,0,1,0,0,0,0,0}]{D}{*****.****}$ in the\n second case.\n\n \\item If $Q$ is of type $D_4^{(1)}$,\n $\\bv^s = \\dynkin[edge length=.6cm,labels={1,2,1,2}]{D}{4}$,\n $\\bw^s = \\dynkin[edge length=.6cm,labels={0,0,0,2}]{D}{4}$, and\n its cyclic permutation.\n\n \\item If $Q$ is of type $E_8^{(1)}$,\n $\\bv^s = \\dynkin[edge\n length=.6cm,labels={4,5,7,10,8,6,4,2}]{E}{8}$,\n $\\bw^s = \\dynkin[edge\n length=.6cm,labels={1,0,0,0,0,0,0,0}]{E}{8}$.\n\\end{itemize}\n\nIn particular, it means that the quiver variety ${\\mathfrak{M}}_0(\\bv^s,\\bw^s)$\nis isomorphic to the intersection of Slodowy slice and the nilpotent\ncone in the formal neighborhood of the origin.\n\nThis result was known before for type $A_n$ and the second case of\n$D_n$. The type $A_n$ case was proved in \\cite[\\S8]{Na-quiver}. See\nalso \\cite{MR2130242}. For the second case of type $D_n$, it was shown\nin \\cite{MR3269179}. In these cases, ${\\mathfrak{M}}_0(\\bv^s,\\bw^s)$ itself is\nisomorphic to the intersection of Slodowy slice and the nilpotent\ncone, not only in the formal neighborhood.\nThe exceptional cases are conjectured, but not shown as far as the\nauthor knows.\n\\end{Example}\n\n\\begin{Remark}\n In cases for the above Example, it is known that corresponding\n Coulomb branches of the quiver gauge theories (which are affine\n Grassmannian slices by \\cite{2016arXiv160403625B}) are\n `next-to-minimal' nilpotent orbits closures in ${\\mathfrak{g}}_{\\mathrm{fin}}$ by\n \\cite{MR3131493}.\n\\end{Remark}\n\n\\begin{Example}[Yamagishi \\cite{2017arXiv170905886Y}]\n Let us consider $\\bv = \\delta+\\alpha_0$, $\\bw = \\Lambda_0$, hence\n ${\\mathfrak{M}}_{\\zeta^\\bullet}(\\delta+\\alpha_0,\\Lambda_0)$. Note that this is\n different from\n ${\\mathfrak{M}}_{\\zeta^\\bullet}(\\delta,\\Lambda_0) =\n \\operatorname{Hilb}^1({\\mathbb C}^2\/\\Gamma) = {\\mathbb C}^2\/\\Gamma$. Let us show\n that ${\\mathfrak{M}}_{\\zeta^\\bullet}(\\delta+\\alpha_0,\\Lambda_0)$ is \\emph{not}\n isomorphic to ${\\mathbb C}^2\/\\Gamma$.\n\n This ${\\mathfrak{M}}_{\\zeta^\\bullet}(\\delta+\\alpha_0,\\Lambda_0)$ appears as the\n union of strata in\n ${\\mathfrak{M}}_{\\zeta^\\bullet}(2\\delta,\\Lambda_0) =\n \\operatorname{Hilb}^2({\\mathbb C}^2\/\\Gamma)$. Thus all possible strata are\n determined already as above: besides the open stratum for\n $\\bv' = \\bv$, we have a single $0$-dimensional stratum except for\n $A_1$, $A_2$, $D_n$, none for $A_1$, $A_2$, three for $D_4$, two for\n $D_n$ ($n > 4$).\n The transversal slice ${\\mathfrak{M}}_0(\\bv^s,\\bw^s)$ is given by the same\n $\\bw^s$ as above, $\\bv^s$ is subtracting\n $\\sum_{i\\in I^0} a_i\\alpha_i$ from the above example. Concretely it\n is\n \\begin{itemize}\n \\item For type $A_n^{(1)}$ with $n > 2$,\n $\\bv^s = \\dynkin[edge\n length=.6cm,labels={0,1,1,1,1,0}]{A}{***.***}$,\n $\\bw^s = \\dynkin[edge\n length=.6cm,labels={0,1,0,0,1,0}]{A}{***.***}$.\n\n \\item For type $D_n^{(1)}$ with $n > 4$,\n $\\bv^s = \\dynkin[edge\n length=.6cm,labels={1,0,0,0,0,0}]{D}{***.***}$,\n $\\bw^s = \\dynkin[edge\n length=.6cm,labels={2,0,0,0,0,0}]{D}{***.***}$ in the first case,\n and\n $\\bv^s = \\dynkin[edge\n length=.6cm,labels={0,0,1,2,2,2,2,1,1}]{D}{*****.****}$,\n $\\bw^s = \\dynkin[edge\n length=.6cm,labels={0,0,0,1,0,0,0,0,0}]{D}{*****.****}$ in the\n second case.\n\n \\item For type $D_4^{(1)}$,\n $\\bv^s = \\dynkin[edge length=.6cm,labels={0,0,0,1}]{D}{4}$,\n $\\bw^s = \\dynkin[edge length=.6cm,labels={0,0,0,2}]{D}{4}$, and\n its cyclic permutation.\n\n \\item For type $E_8^{(1)}$,\n $\\bv^s = \\dynkin[edge\n length=.6cm,labels={2,2,3,4,3,2,1,0}]{E}{8}$,\n $\\bw^s = \\dynkin[edge\n length=.6cm,labels={1,0,0,0,0,0,0,0}]{E}{8}$.\n \\end{itemize}\n\n Let us consider\n \\(\n \\pi_{\\zeta^\\bullet,\\zeta}\\colon {\\mathfrak{M}}_{\\zeta}(\\bv,\\bw)\n \\to {\\mathfrak{M}}_{\\zeta^\\bullet}(\\bv,\\bw).\n \\)\n The domain ${\\mathfrak{M}}_{\\zeta}(\\bv,\\bw)$ is the minimal resolution of\n ${\\mathbb C}^2\/\\Gamma$. The above description of the transversal slice\n immediately conclude the following:\n \\begin{itemize}\n \\item ${\\mathfrak{M}}_{\\zeta^\\bullet}(\\bv,\\bw)$ is obtained from the minimal\n resolution of ${\\mathbb C}^2\/\\Gamma$ by collapsing $\\mathbb P^1$'s whose\n vertices are not connected to the $0$-vertex in the affine Dynkin\n diagram.\n \\end{itemize}\n\n There are no such vertices for $A_1$, $A_2$. For type $A_n$\n ($n > 2$), we collapse $(n-2)$ $\\mathbb P^1$'s corresponding to\n vertices except the leftmost and rightmost. For type $D_n$, we\n collapse $(n-1)$ $\\mathbb P^1$'s except the second one from the\n left, and producing singularities of type $A_1$ and $D_{n-2}$, where\n we understand $D_2 = A_1\\times A_1$, $D_3 = A_3$.\n For type $E_8$, the $7$ $\\mathbb P^1$'s except the rightmost one are\n collapsed.\n\\end{Example}\n\n\\subsection{Euler numbers}\n\nBy \\cref{subsec:strat} we can relate Euler numbers of\n${\\mathfrak{M}}_{\\zeta}(\\bv,\\bw)$, ${\\mathfrak{M}}_{\\zeta^\\bullet}(\\bv',\\bw)$ and\n${\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)$. Let $\\chi(\\ )$ denote the Euler number of\na space. We have\n\\begin{equation*}\n \\begin{split}\n \\chi({\\mathfrak{M}}_{\\zeta}(\\bv,\\bw))\n & = \\sum_{\\bv'} \\chi({\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\\bv',\\bw))\n \\chi(\\mathfrak L_{\\zeta}(\\bv^s,\\bw^s)) \\\\\n & = \\sum_{\\bv'} \\chi({\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\\bv',\\bw))\n \\chi({\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)).\n \\end{split}\n\\end{equation*}\nIt is known \\cite[Cor.~5.5]{Na-quiver} that the central fiber\n$\\mathfrak L_{\\zeta}(\\bv^s,\\bw^s)$ of\n${\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)\\to {\\mathfrak{M}}_0(\\bv^s,\\bw^s)$ is homotopic to\n${\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)$, hence the second equality follows.\n\nFrom $\\chi({\\mathfrak{M}}_{\\zeta}(\\bv,\\bw))$,\n$\\chi({\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s))$, we compute\n$\\chi({\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\\bv',\\bw))$ recursively as\n\\begin{equation}\\label{eq:5}\n \\chi({\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\\bv,\\bw))\n = \\chi({\\mathfrak{M}}_\\zeta(\\bv,\\bw)) - \\sum_{\\bv'\\neq \\bv}\n \\chi({\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\\bv',\\bw))\n \\chi({\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)).\n\\end{equation}\nHere we use $\\bv^s = 0$, hence ${\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)$ is a point\nfor $\\bv'=\\bv$.\n\nWe take the generating function of Euler numbers as\n\\begin{equation}\\label{eq:8}\n \\sum_{\\bv} \\chi({\\mathfrak{M}}_{\\zeta}(\\bv,\\bw)) e^{-\\bv}\n = \\sum_{\\bv'} \\chi({\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\\bv',\\bw)) e^{-\\bv'}\n \\sum_{\\bv^s}\n \\chi({\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw-\\bv'|_{I^0})) e^{-\\bv^s},\n\\end{equation}\nwhere\n$\\bw - \\bv'|_{I^0} = \\bw^s = \\sum_{i\\in I^0} \\langle \\bw -\n\\bv',\\alpha_i^\\vee\\rangle \\Lambda_i$.\n\nWe claim\n\\begin{equation}\\label{eq:1}\n \\left.\\sum_{\\bv^s}\n \\chi({\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)) \n \\prod_{i\\in I^0} e^{\\bw^s_i \\Lambda_i\n - \\bv^s_i \\alpha_i}\n \\right|_{\n e^{-\\alpha_i} = \\exp(\\frac{2\\pi i}{h^\\vee+1})\n }\n = 1.\n\\end{equation}\n\\begin{NB}\n $\\mathop{\\text{\\rm dim}}\\nolimits {\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s) = \\langle 2\\bw^s-\\bv^s, \\bv^s\\rangle$.\n\\end{NB}%\n\nWe take $\\bw = \\Lambda_0$ and switch to the convention\n\\eqref{eq:3}. Then the left hand side of \\eqref{eq:8} is $Z$ in\nIntroduction. By \\eqref{eq:6}, \\eqref{eq:1} implies that\n\\begin{NB}\n\\begin{multline*}\n \\sum_{m, \\vec{m}} \\chi({\\mathfrak{M}}_{\\zeta}(m\\delta+\\sum_{i\\in I^0} m_i\\alpha_i,\n \\Lambda_0)) \n e^{-m\\delta} \\prod_{i\\in I^0} e^{-m_i\\alpha_i}\n\\\\\n = \\sum_{m,\\vec{m}'} \\chi({\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\n m\\delta+\\sum_{i\\in I^0} m'_i\\alpha_i,\\Lambda_0)) \n e^{-m\\delta} \\prod_{i\\in I^0} e^{-m_i\\alpha_i}\n \\sum_{\\vec{m''}}\n \\chi({\\mathfrak{M}}_{\\zeta}(\\sum_{i\\in I^0} m''_i\\alpha_i,\\bw^s)) \n \\prod_{i\\in I^0}e^{-m''_i\\alpha_i}.\n\\end{multline*}\n\\end{NB}\n\\begin{equation*}\n \\left.Z\\right|_{\n e^{-\\alpha_i} = \\exp(\\frac{2\\pi i}{h^\\vee+1})\n }\n = \\sum_{m,\\vec{m}'} \\chi({\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}(\n m\\delta+\\sum_{i\\in I^0} m_i'\\alpha_i,\\Lambda_0)) \n e^{-m\\delta}.\n\\end{equation*}\nBy \\eqref{eq:7} we consider\n${\\mathfrak{M}}_{\\zeta^\\bullet}^{\\mathrm{s}}( m\\delta+\\sum_{i\\in I^0}\nm_i'\\alpha_i,\\Lambda_0)$ as a stratum of\n${\\mathfrak{M}}_{\\zeta^\\bullet}(m\\delta,\\Lambda_0) =\n\\operatorname{Hilb}^m({\\mathbb C}^2\/\\Gamma)$. Therefore the right hand side\nof the above is the generating function of Euler numbers of\n$\\operatorname{Hilb}^m({\\mathbb C}^2\/\\Gamma)$.\nThus we have proved \\cref{thm:1}.\n\nRecall that\n$\\bigoplus_{\\bv^s} H_*(\\mathfrak L_{\\zeta}(\\bv^s,\\bw^s)) \\cong\n\\bigoplus_{\\bv^s} K(\\mathfrak L_{\\zeta}(\\bv^s,\\bw^s))\\otimes_{{\\mathbb Z}}{\\mathbb C}$\nis the so-called standard module of the quantum loop algebra\n$\\mathbf U_q(\\mathbf L{\\mathfrak{g}}_{\\mathrm{fin}})$, specialized at $q=1$\n\\cite{Na-qaff}.\nThe above \\eqref{eq:1} means that its character (as a\n${\\mathfrak{g}}_{\\mathrm{fin}}$-module), specialized at\n$e^{-\\alpha_i} = \\exp(\\frac{2\\pi i}{h^\\vee+1})$ is equal to $1$.\nAs we mentioned in Introduction, this specialization is the quantum\ndimension. Thus \\eqref{eq:1} follows from \\cref{thm:2}.\n\nIn order to prove \\eqref{eq:1}, we may assume $\\bw^s$ is a fundamental\nweight: there is a torus action on framing vector spaces, and the\ninduced action on ${\\mathfrak{M}}_{\\zeta}(\\bv^s,\\bw^s)$. Let us suppose\n$\\bw^s = \\Lambda_i+\\Lambda_j$ for $i,j\\in I^0$ for simplicity. Then\nthe torus fixed point set is\n\\begin{equation*}\n \\bigsqcup_{\\bv^1+\\bv^2=\\bv^s} {\\mathfrak{M}}_{\\zeta}(\\bv^1,\\Lambda_i)\n \\times {\\mathfrak{M}}_{\\zeta}(\\bv^2,\\Lambda_j).\n\\end{equation*}\nAs Euler number is equal to the sum of Euler numbers of fixed points\nwith respect to a torus action, \\eqref{eq:1} for $\\bw^s$ follows from\n\\eqref{eq:1} for $\\Lambda_i$ and $\\Lambda_j$.\nThis result is compatible with what we mentioned in Introduction:\nstandard modules are tensor products of $l$-fundamental modules\n\\cite{VV-std}, as $l$-fundamental modules correspond to the case\n$\\bw^s$ is a fundamental weight.\n\nA standard module depends also on spectral parameter, which is\nspecialization homomorphism\n$K_{\\prod_{i\\in I^0} \\operatorname{GL}(\\bw^s_i)}(\\mathrm{pt})\\to {\\mathbb C}$. But the\nrestriction of a standard module to $\\mathbf U_q({\\mathfrak{g}}_{\\mathrm{fin}})$ is\nindependent of the spectral parameter. Hence the spectral parameter is\nnot relevant for \\cref{thm:2}.\n\n\\begin{Remark}\n In order to prove \\cref{thm:1}, we need to check \\eqref{eq:1} only\n when $\\bw^s$ is contained in the root lattice. But $\\Lambda_i$ does\n not satisfy this condition in general, hence the above reduction\n cannot be performed among $\\bw^s$ in the root lattice.\n\\end{Remark}\n\n\\section{Quantum dimensions of standard modules}\\label{sec:2}\n\n\\subsection{Quantum dimension}\n\nLet $V$ be a finite dimensional representation of ${\\mathfrak{g}}_{\\mathrm{fin}}$. The\nspecialized character\n\\begin{equation*}\n \\left.\\operatorname{ch} V\\right|_{e^{-\\alpha_i}=\\exp(\\frac{2\\pi i}{h^\\vee+1})}\n\\end{equation*}\nis called the \\emph{quantum dimension} of $V$, denoted by $\\mathop{\\text{\\rm dim}}\\nolimits_q V$,\nand was introduced by Andersen~\\cite[Def.~3.1]{MR1182414} (see also\nParshall-Wang~\\cite{MR1247547}), where\n$\\zeta = \\exp(\\frac{2\\pi i}{2(h^\\vee+1)})$. Note $2\\rho$ in\n$K_{2\\rho} = q^{2\\rho}$ in \\cite{MR1182414} should be understood as an\nelement in the dual of the weight lattice by\n$\\langle 2\\rho,\\lambda\\rangle = (2\\rho,\\lambda)\/(\\alpha_0,\\alpha_0)$\nfor a weight $\\lambda$. Here $\\alpha_0$ is a short root. See\n\\cite[Lemma~1.1]{MR1247547}.\nIt was assumed that $\\zeta$ is a primitive $\\ell$-th root of unity\nwith odd $\\ell$ in \\cite{MR1182414}, the definition still makes sense\nfor our choice. By the Weyl character formula, we have (see\n\\cite[(3.2)]{MR1182414}, \\cite[Th.~1.3]{MR1247547})\n\\begin{equation*}\n \\left.\\operatorname{ch} V\\right|_{e^{-\\alpha_i}=\\exp(\\frac{2\\pi i}{h^\\vee+1})}\n = \\mathop{\\text{\\rm dim}}\\nolimits_q V\n = \\prod_{\\alpha\\in\\Delta^+_{\\mathrm{fin}}} \\frac{\n \\zeta^{d_\\alpha\\langle\\lambda+\\rho,\\alpha^\\vee\\rangle} -\n \\zeta^{-d_\\alpha\\langle\\lambda+\\rho,\\alpha^\\vee\\rangle}}{\n \\zeta^{d_\\alpha\\langle\\rho,\\alpha^\\vee\\rangle} -\n \\zeta^{-d_\\alpha\\langle\\rho,\\alpha^\\vee\\rangle}},\n\\end{equation*}\nwhere $\\lambda$ is the highest weight of an irreducible representation\n$V = V(\\lambda)$, $\\Delta^+_{\\mathrm{fin}}$ is the set of positive roots of\n${\\mathfrak{g}}_{\\mathrm{fin}}$, $d_\\alpha\\in \\{1,2,3\\}$ is the square length of $\\alpha$\ndivided by the length of a short root, and $\\rho$ is the half sum of\npositive roots. (Since we are only considering type $ADE$, we have\n$d_\\alpha = 1$ for any $\\alpha$.)\nWe have\n$d_\\alpha\\langle\\rho,\\alpha^\\vee\\rangle\\le \\langle\\rho,\\theta\\rangle =\nh^\\vee$, hence the denominator does not vanish.\n\nLet us introduce $\\zeta$-integers by\n\\begin{equation*}\n [n]_{\\zeta} \\overset{\\operatorname{\\scriptstyle def.}}{=} \\frac{\\zeta^{n} - \\zeta^{-n}}\n {\\zeta - \\zeta^{-1}},\n\\end{equation*}\nso that\n\\begin{equation}\\label{eq:4}\n \\left.\\operatorname{ch} V\\right|_{e^{-\\alpha_i}=\\exp(\\frac{2\\pi i}{h^\\vee+1})}\n = \\prod_{\\alpha\\in\\Delta^+}\n \\frac{[d_\\alpha\\langle\\lambda+\\rho,\\alpha^\\vee\\rangle]_{\\zeta}}\n {[d_\\alpha\\langle\\rho,\\alpha^\\vee\\rangle]_{\\zeta}}.\n\\end{equation}\nWe have\n\\begin{equation*}\n [2(h^\\vee+1) + k]_{\\zeta} = [k]_{\\zeta}\n\\end{equation*}\nas $\\zeta^{2(h^\\vee+1)} = 1$, as well as\n\\begin{equation}\\label{eq:2}\n [h^\\vee+1-k]_{\\zeta} = [k]_{\\zeta}\n\\end{equation}\nas $\\zeta^{h^\\vee+1-k} \\zeta^{k} = \\zeta^{h^\\vee+1} = -1$. In\nparticular, we have $[h^\\vee+1]_{\\zeta} = 0$. The former is usual\nanalogy between roots of unity and characteristic $2(h^\\vee+1)$. The\nlatter is a new feature at an even root of unity.\n\n\\subsection{type $A$}\n\nConsider type $A_{n-1}$, i.e., ${\\mathfrak{g}}_{\\mathrm{fin}} = \\algsl(n,{\\mathbb C})$. We have\n$h^\\vee = n$.\n\nIn type $A_{n-1}$, it is known that the $k$th $l$-fundamental module\nis the $k$th fundamental representation of $\\operatorname{\\rm SL}(n)$.\n\nThe quantum dimension of the $k$th fundamental representation\n$V(\\Lambda_k)$ is\n\\begin{equation*}\n \\mathop{\\text{\\rm dim}}\\nolimits_q V(\\Lambda_k)\n = \\begin{bmatrix}\n n \\\\ k\n \\end{bmatrix}_{\\zeta}\n = \\frac{[n]_{\\zeta}[n-1]_{\\zeta}\\dots [n-k+1]_{\\zeta}}\n {[k]_{\\zeta} [k-1]_{\\zeta} \\dots [1]_{\\zeta}}.\n\\end{equation*}\n\\begin{NB}\nIn fact, the denominator of \\eqref{eq:4} is equal to\n\\begin{equation*}\n \\prod_{1\\le i < j\\le n} [j-i]_\\zeta,\n\\end{equation*}\nas positive roots $\\alpha$ are\n$\\alpha_{i+1} + \\dots + \\alpha_j$.\n\nOn the other hand, from the numerator, we have\n\\begin{equation*}\n \\prod_{1\\le i < j\\le n} [\\lambda_i - \\lambda_j + j-i]_\\zeta,\n\\end{equation*}\nwhere $\\lambda$ is regarded as a partition $(1^k)$.\nSince $\\lambda_i - \\lambda_j$ is $1$ if $i\\le k$, $j>k$ and $0$\notherwise, we have\n\\begin{equation*}\n \\prod_{1\\le i 2$. We have\n $\\bw^s = \\Lambda_2 + \\Lambda_{n-1}$. The standard module is\n ${\\textstyle \\bigwedge}^2 {\\mathbb C}^{n+1}\\otimes {\\textstyle \\bigwedge}^{n-1} {\\mathbb C}^{n+1}$. It contains the\n trivial representation with multiplicity $1$.\n\n Consider type $D_n$. For $\\bw^s = 2\\Lambda_1$, the standard module\n is the tensor square of the vector representation, which is\n $\\operatorname{Sym}^2{\\mathbb C}^{2n}\\oplus {\\textstyle \\bigwedge}^2{\\mathbb C}^{2n}$, and the\n former is the sum of the irreducible representation of the highest\n weight $2\\Lambda_1$ and the trivial representation. (See e.g.,\n \\cite[Exercise~19.21]{MR1153249}. Therefore the multiplicity of the\n trivial representation is $1$.\n\\end{NB}%\n\n\\subsection{Counter-example}\n\nConsider $A_1^{(1)}$ with $n=4$, i.e.,\n$\\operatorname{Hilb}^4({\\mathbb C}^2\/({\\mathbb Z}\/2)) =\n{\\mathfrak{M}}_{\\zeta^\\bullet}(4\\delta,\\Lambda_0)$. We take\n$\\bv' = \\dynkin[arrows=false,edge length=.6cm,labels={4,2}]{B}{o*}$.\nThen $\\bv^s = 2$, $\\bw^s = 4$. Hence the transversal slice is\n${\\mathfrak{M}}_0(\\bv^s,\\bw^s)$, which is the nilpotent orbit in\n$\\mathfrak{sl}(4)$ of type $(2^2)$. It is well known to be a\nnon-rational homology manifold. For example, we can argue as\nfollows. By \\cite[Th.~15.1.1]{Na-qaff}\n$\\mathop{\\text{\\rm dim}}\\nolimits i_0^! \\mathrm{IC}({\\mathfrak{M}}_0(\\bv^s,\\bw^s))$ gives the multiplicities\nof the trivial representation in the standard module for $\\bw^s$. The\nlatter is $4$th tensor power of the natural $2$-dimensional\nrepresentation of $\\algsl(2)$. The multiplicity of the trivial\nrepresentation is $2$.\n\\begin{NB}\n $(y+y^{-1})^4 = y^4 + 4y^2 + 6 + 4y^{-2} + y^{-4}\n = (y^4 + y^2 + 1 + y^{-2} + y^{-4}) + 3(y^2 + 1 + y^{-2}) + 2$.\n\\end{NB}%\nHence ${\\mathfrak{M}}_0(\\bv^s,\\bw^s)$ is \\emph{not} a rational homology manifold,\nand ${\\mathfrak{M}}_{\\zeta^\\bullet}(4\\delta,\\Lambda_0)$ neither.\n\nAlternatively we argue as follows.\nWe realize this variety as the Coulomb branch of a quiver\ngauge gauge theory of type $A_3$ with\n$\\bv'' = \\dynkin[edge length=.6cm,labels={1,2,1}]{A}{3}$,\n$\\bw'' = \\dynkin[edge length=.6cm,labels={0,2,0}]{A}{3}$. It is the\nclosure of a $\\operatorname{\\rm SL}(4)[[z]]$-orbit in the affine Grassmannian for\n$\\operatorname{\\rm SL}(4)$, and hence the intersection cohomology is known by geometric\nSatake correspondence. We have\n\\begin{equation*}\n \\mathop{\\text{\\rm dim}}\\nolimits i^!_0\\mathrm{IC}(\\mathcal M_C(\\bv'',\\bw'')) = \\mathop{\\text{\\rm dim}}\\nolimits V_0(2\\Lambda_2),\n\\end{equation*}\nwhere $i_0$ is the embedding of the identity element to the affine\nGrassmannian. The right hand side is the $0$-weight space in the\nrepresentation of $\\operatorname{\\rm SL}(4)$ with the highest weight $2\\Lambda_2$. This\nweight space is $2$-dimensional.\n\nIn general, if ${\\mathfrak{M}}_0(\\bv^s,\\bw^s)$ is happen to be the closure of a\n$G[[z]]$-orbit in an affine Grassmannian $\\mathrm{Gr}_G = G((z))\/G[[z]]$\n(e.g., it is true in type $A$ by \\cite{MR1968260} or the combination\nof \\cite{2016arXiv160602002N} and \\cite{2016arXiv160403625B}), it is a\nrational homology manifold if and only if weight spaces of the\ncorresponding representation are all $1$-dimensional.\\footnote{The\n author thanks Dinakar Muthiah for this remark.}\n\n\\bibliographystyle{myamsalpha}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\nIn any information system which requires user interaction, being able to understand the user is a crucial requirement, which is often tackled by limiting the user input (e.g., presenting predefined forms with fixed options). The more freedom you provide the user with, the more difficult interpretation the computer has to do to achieve a useful interaction. In such a context, being capable of disambiguating the input words (i.e., associating each word with its proper meaning in a given context) is the starting point of any interpretation process done by the computer. \n\nUsually, such a disambiguation process is tackled from a Natural Language Processing (NLP) perspective~\\cite{Navigli2009}, assuming rich linguistic information, such as Part of Speech (POS), dependencies between words, etc., which is very useful to perform the task. However, due to the world wide use of Web search engines, users are very used to keyword interfaces and they still express their needs in such terms. In this scenario, although there are some studies which point out that \\emph{keyword search queries} (aka., \\emph{Web search queries}) exhibit their own language structure~\\cite{barr2008linguistic,pinter2016syntactic,roy2016syntactic}, we still need methods to disambiguate the meanings of the words which do not need such an information as it might not be available. \n\n\nRecent advances in NLP have focused on the development of different embedding models~\\cite{bengio2003neural,mikolov2013distributed,le2014distributed,mancini2017embedding,camacho2016nasari,pennington2014glove}, which are a set of language modeling and feature learning techniques where elements from a vocabulary are mapped to a vectorial space capturing their \\emph{distributional semantics}~\\cite{sahlgren2008distributional}.\nWhile there are different methods to build word embeddings, the latest (and most successful) word embedding techniques rely on neural network architectures~\\cite{bengio2003neural,mikolov2013distributed,le2014distributed,pennington2014glove}.\nTheir usage as the underlying input representation has boosted the performance of different NLP tasks~\\cite{socher2013parsing,socher2013recursive}. However, in the context of disambiguation tasks, one of the main limitations of word embeddings is that the possible meanings of a word are combined into a single representation, i.e., a single vector in the semantic space. Such a limitation can be avoided by representing individual meanings of words as distinct vectors in the space (e.g., sense2vec~\\cite{mancini2017embedding}). However, there are scenarios where we do not know all the different senses at training time (e.g., open domain scenarios where we cannot find all the possible meanings in a sense catalog), and, even if we would know them, we would require to have annotated data (which might be unavailable or expensive to obtain). Besides, we would need to train a model for each new scenario or new meaning that would be added to our catalog. Thus, we need a disambiguation method able to relate words and their senses in a flexible and general way (i.e., independently of the domain we are working in) by exploiting the available resources.\n\nIn this paper, we propose a keyword disambiguation method which is based on the semantic relatedness (the degree in which two objects are related by any kind of semantic relationship~\\cite{budan2006evaluating}) between words, taking advantage of the semantic information captured by word embeddings. Our proposal makes possible to measure the relatedness not only among plain words but also among senses of words (which, in a Semantic Web context, can be expressed as ontological terms), and it is able to work independently of the resources used, i.e., the sense inventory whose meanings we want to map to and the word embedding model used as input.\n\nFor this purpose, we build on the work by Gracia~and~Mena on semantic relatedness~\\cite{wise08} and disambiguation~\\cite{gracia2009multiontology}. These works exploited the information about word co-occurrence frequencies provided by existing Web search engines. We evolve and adapt them to improve their performance using different kinds of embeddings (both at word and sentence level).\nThe main benefit of such an adaptation is two-fold:~1) we exploit the semantics captured by embeddings which goes further than co-occurrence of terms, and~2) we decouple the proposal from any Web search engine, being able to use off-the-shelf models trained by third parties for our purposes. This has an important side-effect: our measure can be easily adapted to any domain which we have a document corpus from. Indeed, this adaptation would require a training step, but it would be unsupervised and the only data required would be the corpus of documents itself. \n\nTo evaluate our approach, we have carried out a thorough experimentation in the context of Word Sense Disambiguation (WSD), where we have used different pre-trained word embeddings publicly available on the Web, and WordNet\\footnote{\\label{wordnetfootnote}\\url{https:\/\/wordnet.princeton.edu\/}} as sense repository. Our measure improves the performance obtained in~\\cite{wise08}, and achieves state of the art WSD values without the need of specific training for a specific sense inventory. This is especially relevant, for example, for systems based on keyword input and\/or which have to work with dynamically selected ontologies~\\cite{bobed2016querygen} or even with ontologies extracted directly from the Web~\\cite{Halevy15}. All the experimental data and evaluation results are available online\\footnote{\\url{https:\/\/bit.ly\/2lqCzop}}.\n\nThe rest of the paper is structured as follows. Section~\\ref{sec:relatedwork} discusses related works. In Section~\\ref{sec:approachdescription} we describe our semantic relatedness measure approximation, in Section~\\ref{sec:disambalgorithm} we present the disambiguation algorithm that we use, and Section~\\ref{sec:experimentaleval} summarizes our experimental results. Finally, our conclusions and future work appear in Section~\\ref{sec:conclusionsandfuturework}.\n\n\\section{Related Work}\n\\label{sec:relatedwork}\nSemantic relatedness is the degree in which two objects are related by any kind of semantic relationship~\\cite{budan2006evaluating} and lies at the core of many applications in NLP (such as WSD, Information Retrieval, Natural Language Understanding, or Entity Recognition). The term is often confused with semantic similarity, which measures the degree in which two objects are similar or equivalent. For example, \"car\" is similar to \"bus\", but is also related to \"road\" and \"driving\". It has received a great research interest and different types of methods have been developed: it can be statistically estimated (e.g. co-occurrence-based methods~\\cite{landauer1998introduction}) and learned (e.g., distributional measures that estimate semantic relatedness between terms using a multidimensional space model to correlate words and textual contexts~\\cite{mohammad2012distributional}); or it can be computed using a taxonomy or a graph (e.g., ontologies) to define the distance between terms or concepts~\\cite{pirro2012reword}. Indeed, most methods rely on particular lexical resources (dictionaries, thesauri, or well structured taxonomies such as WordNet\\textsuperscript{\\ref{wordnetfootnote}}).\n\nRegarding disambiguation, WSD methods can be classified into four conventional approaches: supervised~\\cite{vial2018improving}, unsupervised~\\cite{correa2018word}, semi-supervised~\\cite{yuan2016semi}, and knowledge-based methods~\\cite{chaplot2018knowledge}. For example, in a way similar to us, in the SemEval 2015 All-Words Sense Disambiguation and Entity Linking task\\footnote{\\label{semeval2015footnote}\\url{http:\/\/alt.qcri.org\/semeval2015\/task13\/}}, the majority of the approaches (LIMSI, SUDOKU, EBL-Hope, etc.) that best performed in WSD relied on the combination of unsupervised learning of semantic information from the content of a corpus (such as SemCor) and\/or on lexical resources as sense inventories (such as WordNet or BabelNet) to disambiguate the sense of words in natural language sentences. However, to our knowledge, no previous works (excepting those of Gracia~and~Mena~\\cite{wise08,gracia2009multiontology}) have studied specific methods for the disambiguation of words in keyword-based inputs, where the linguistic information is scarce.\n\nRegarding the resources we use in our approach, word embeddings represent words in a low-dimensional continuous space and they are used to capture syntactic and semantic information from massive amounts of textual content. In recent years, they have gained great popularity due to this ability and many NLP applications have taken advantage of the potential of these distributional models. Bengio et al.~\\cite{bengio2003neural} preceded a wide number of current language model techniques and several authors have proposed their own approaches~\\cite{le2014distributed,pennington2014glove,camacho2016nasari} to construct word embeddings vectors where \\emph{word2vec}~\\cite{mikolov2013distributed} is the most widely used.\n\nDespite their advantages, one of the main limitations of word embeddings is that possible meanings of a word are conflated into a single representation. Sense embeddings (e.g., sense2vec~\\cite{mancini2017embedding}) are proposed as a solution to this problem: individual meanings of words are represented as distinct vectors in the space. These approaches are classified in two categories by how they model meaning and where they obtain it from~\\cite{camacho2018word}: 1) unsupervised models which learn word senses from text corpora (by inducing different senses of a word, analyzing its contextual semantics in a text corpus and representing each sense based on the statistical knowledge derived from the corpus), and 2) knowledge-based methods which exploit sense inventories of lexical resources for representing meanings (such as WordNet\\textsuperscript{\\ref{wordnetfootnote}}, Wikipedia\\footnote{\\label{wikifootnote}\\url{https:\/\/www.wikipedia.org\/}}, BabelNet\\footnote{\\label{bnfootnote}\\url{https:\/\/babelnet.org\/}}, etc.). We can also find models that not only provide representations of words, but also the senses of the words in a joint embedded space. This is the case of NASARI vectors~\\cite{camacho2016nasari} which not only provide accurate representation of word senses in different languages, but they also include both concepts and named entities, all included in a single unified semantic space. \nHowever, in the first case (i.e., unsupervised models), we cannot target a particular sense inventory or ontology to perform the disambiguation, not having control for example about the concept detail\/granularity, and, besides, the detected senses might not be aligned to any particular human-readable structure; in the second case, we need to know all the senses at training time, not being able to adapt to new scenarios (e.g., addition\/deletion of senses in the directory, evolving ontologies, etc.). Thus, the sweet spot would be neither to require re-training nor newly labelled data, while being capable to perform the disambiguation against any sense repository.\n\n\n\n\nAlthough sense embeddings capture and represent information about meanings and can be used to calculate the sense that a word has in a specific context, word embeddings have also been shown to have good performance in disambiguation tasks~\\cite{iacobacci2016embeddings}. Therefore, we wanted to explore how we could push further the usage of word embeddings for keyword disambiguation. Working at word level (as starting point) allows us to use a semantic relatedness measure between terms and reuse available resources, without needing to train explicitly new word embeddings neither for a specific task nor for newly added possible senses (i.e., adapting to any given sense dictionary or ontology). We have taken as baseline the works presented in~\\cite{wise08,gracia2009multiontology}. In~\\cite{gracia2009multiontology}, the authors provide a keyword disambiguation algorithm that uses the semantic relatedness measure defined in~\\cite{wise08} to find the appropriate sense for keywords. The authors focused on a method that exploits the Web as a source of knowledge, and a transformation of the Normalized Google Distance (NGD)~\\cite{cilibrasi2007google} into a mixed way of relatedness measure (between ontology terms and plain words). \nWe propose to, on the one hand, substitute this distance with a word embedding based one to take advantage of the semantics captured by embeddings, improving the performance regarding using just co-occurrence of terms; and, on the other hand, explore modifications of their algorithm to improve their disambiguation capabilities.\n\nFinally, as pointed out by Lastra-D\\'iaz~et~al.~\\cite{lastra2019reproducible}, the embeddings that behave the best for disambiguation purposes are those which capture not only distributional semantics of texts, but also structural information about the possible meanings. We aim at achieving this disambiguation performance in a more flexible way, decoupling linguistic surface from the actual sense catalog (i.e., ontology) in order to adapt to new (i.e., unknown at training time) possible meanings, and being capable to apply it to keyword inputs, where the linguistic information is scarce.\n\n\\section{Relatedness Measure based on Word Embeddings}\n\\label{sec:approachdescription}\nWord embeddings can be used out-of-the-box to compute relatedness between words. However, they do not suffice in situations in which a relatedness has to be computed between senses (e.g., ontology terms in a Semantic Web context) or between senses and words. To that end, we ground on a previously defined relatedness measure between senses proposed by Gracia~and~Mena~\\cite{wise08}.\nThe authors proposed a method to compute the semantic relatedness between ontology terms (which we can see as individual senses), and an extension to calculate it between plain words and terms.\nTheir proposal was built on the notion of the \\emph{ontological context} of a term, which is constructed combining the synonyms and the hypernyms of the ontological term (or sense). Given an ontological term ${t}$, they defined its \\emph{ontological context} (denoted by ${OC(t)}$) as the minimum set of other ontological terms that belong to its semantic description, locating the term in the ontology and characterizing its meaning. For example, in the WordNet taxonomy, the class ``Java'' (in the sense of ``an Indonesian island''), is well characterized and distinguished from other senses by considering its direct hypernym ``Island'' (see Figure~\\ref{fig:Java}).\n\n\\begin{figure}\n\\includegraphics[width=8cm]{sense_example.png}\n\\caption{Example of the semantic description of the term \"Java\" in WordNet.}\n\\label{fig:Java}\n\\end{figure}\n\n\nThen, given two ontological terms $a$ and $b$, their relatedness measure is computed as: \n\n\\begin{align}\n\\label{eq:completeeq}\n\\begin{split}\n rel(a,b)=w_0rel_0(a,b)+w_1rel_1(a,b), \\\\\n (w_0 \\geq 0, w1 \\geq 0, w_0+w_1=1) \n\\end{split}\n\\end{align}\n\n\\noindent\nwith ${rel_0(a,b)}$ and ${rel_1(a,b)}$ computed as follows: \n\\begin{align}\n\\label{eq:level0}\n\\begin{split}\nrel_0(a,b)=\\frac{\\sum_{i, j}rel_{w}({syn_a}_i,{syn_b}_j)}{|Syn(a)||Syn(b)|}, \\\\\n(i=1..|Syn(a)|, j=1..|Syn(b)|)\n\\end{split}\n\\end{align}\n\\begin{align}\n\\label{eq:level1}\n\\begin{split}\nrel_1(a,b)=\\frac{\\sum_{i, j}rel_0({oc_a}_i,{oc_b}_j)}{|OC(a)||OC(b)|}, \\\\\n(i=1..|OC(a)|, j=1..|OC(b)|)\n\\end{split} \n\\end{align}\n\n\\noindent where ${rel_{w}}$ refers to the relatedness between words (as it will be defined later on in equations 7 and 9); ${Syn(a)=\\{{syn_a}_1, {syn_a}_2, ...\\}}$ and ${Syn(b)=\\{{syn_b}_1, {syn_b}_2, ...\\}}$ are the set of synonyms (equivalent labels, including the term label) of ontological terms \\textit{a} and \\textit{b}; ${OC(a)=}\\\\{\\{{oc_a}_1, {oc_a}_2, ...\\}}$ and ${OC(b)=\\{{oc_b}_1, {oc_b}_2, ...\\}}$ are the terms of their ontological context\\footnote{Notice that ${|Syn(x)| \\geq 1}$ and ${|OC(x)| \\geq 0}$.}. Each \\textit{a} and \\textit{b} is characterized by taking into account two levels of their semantic description: \\textit{Level 0)} the term label and its synonyms (Equation~\\ref{eq:level0}), and \\textit{Level 1)} its ontological context (Equation~\\ref{eq:level1}). $w_0$ and $w_1$ are used to weight these levels\\footnote{We set these values as $w_0 = w_1 = 0.5$ as indicated in~\\cite{wise08}.}.\n\nThis measure can be also applied between an ontology term ${t}$ and a plain word ${w}$ which provides us with a value which indicates the relatedness degree between a sense and a word. So, in that case, the previous equations are computed as follows:\n\\begin{align}\n\\begin{split}\nrel(t,w)=w_0rel_0(t,w)+w_1rel_1(t,w), \\\\ \n(w_0 \\geq 0, w1 \\geq 0, w_0+w_1=1)\n\\end{split} \\label{eq:termWord}\n\\end{align}\n\\begin{equation}\nrel_0(t,w)=\\frac{\\sum_{i,j}rel_{w}({syn_t}_i,w)}{|Syn(t)|}, \\\\\n(i=1..|Syn(t)|)\n\\end{equation}\n\\begin{equation}\nrel_1(t,w)=\\frac{\\sum_{i,j}rel_0({oc_t}_i,w)}{|OC(t)|}, \\\\\n(i=1..|OC(t)|)\n\\end{equation}\n\nHere, ${rel_{w}}$ is the distance that the authors used in~\\cite{wise08} to measure how two plain words are related. They proposed a generalization of the Cilibrasi and Vit\\'anyi's Normalized Google Distance \\textit{NGD(x,y)}~\\cite{cilibrasi2007google} to use any Web search engine as source of frequencies. This generalization is called Normalized Web Distance \\textit{NWD(x,y)}, whose smaller values represent greater semantic relation between words. Although most of NWD values fall between 0 and 1, it ranges from 0 to ${\\infty}$. Therefore, to obtain a proper relatedness measure in the range [0, 1] that increases inversely to distance, they proposed the following transformation:\n\\begin{equation}\n\\label{eq:relweb}\nrel_w(x, y)=relWeb(x, y)=e^{-2NWD(x,y)}\n\\end{equation}\n\n\n\nTo explore the use of emerging word-embedding techniques in this context and compare them with those based on search engines, we propose to exploit the semantic capabilities of word embeddings in this formulation and substitute the ${relWeb(x,y)}$ measure. We could use the cosine similarity distance between the embedding vectors of the words, i.e., using the following equation:\n\\begin{equation}\nsim(x,y)=cos(\\theta)=\\frac{X1 \\bullet X2}{||X1||\\cdot||X2||}\n\\end{equation}\n\\noindent where ${x}$ and ${y}$ are plain words, ${X1}$ and ${X2}$ their correspondent word embedding vectors, and ${\\theta}$ the angle between them. However, ${sim(x,y)}$ ranges in [-1, 1], so, in order to obtain a \ndistance in the range [0, 1] (so that we can substitute Equation~\\ref{eq:relweb} directly in Equation~\\ref{eq:level0}), we propose to use the \\textit{angular distance} instead, which is computed as follows:\n\\begin{equation}\n\\label{eq:arccos}\nrel_w(x, y)=ang. distance(x,y)=1-\\frac{arccos(sim(x,y))}{\\pi}\n\\end{equation}\n\n\\noindent So, in Equation~\\ref{eq:level0}, we use Equation~\\ref{eq:arccos} as ${rel_{w}}$ distance instead of Equation~\\ref{eq:relweb}.\nWe use this distance to compute the semantic relatedness between words, between ontology terms (or senses), or between ontology terms and words, obtaining a value between 0 and 1. For those cases in which the label of the ontological term is multi-word, we just compute the centroid of the set of words that form the label. While, at first, it might seem that we limit the coverage of the measure proposed in~\\cite{wise08} (it built on the results of Web search engines, which potentially cover any domain), we have to bear in mind the plethora of word embedding models directly available in the Web, as well as the possibility of using our own corpus of documents to fine tuning our measure for a particular domain (which is easier to have, rather than crawling the whole WWW).\n\n\\section{Disambiguation Algorithm}\n\\label{sec:disambalgorithm}\nWe ground our keyword disambiguation proposal on the disambiguation algorithm defined in~\\cite{gracia2009multiontology}, using the adapted semantic relatedness measure proposed in the previous section. This algorithm is based on the hypothesis that the most significant words in the disambiguation context are the most highly related to the word to disambiguate; such words conform the \\emph{active context} of the word being disambiguated.\n\nAs an overview, once the \\emph{active context} of each input keyword has been calculated, the algorithm performs three main steps:~1)~obtaining the semantic relatedness between the active context of a keyword and its possible senses,~2) calculating the overlap between the words in the active context and the semantic descriptions (i.e., \\emph{ontological context}) of the\npossible senses of the keyword to disambiguate, and 3) re-ranking the possible senses according to their frequency of use (only when such information is available for the sense inventory selected\\footnote{If we do not have such information, we assume that all senses are equally likely.}).\nApart from using the updated $rel_w$ measure to select the \\emph{active contexts}, we propose to modify the second step of this algorithm in order to study the influence of different approaches which exploit the semantic information captured by different word embeddings.\nIn the following subsections, we first detail the original algorithm which we base our proposal on, and, then, we describe the modifications that we propose to improve its performance using word embeddings.\n\n\\subsection{Background: Algorithm Description}\n\\label{sec:algorithmdescription}\nFirst of all, let us formally introduce the notion of \\emph{active context}. Let~${k}$ be an element of an input sequence of words~${\\mathbb{S}}$ with an intended meaning, ${\\mathbb{K} \\subseteq \\mathbb{S}}$ be the set of all possible keywords in the input, ${C \\subseteq \\mathbb{K}}$ the set of keywords of the disambiguation context (i.e., the complete disambiguation window considered, e.g., the sentence where the keyword appear), and ${k_d \\in \\mathbb{K}}$ the target keyword to disambiguate. \nThus:\n\\begin{definition}\nGiven a context $\\textit{C} \\in \\mathbb{K}$, and a word to disambiguate ${k_d} \\in \\mathbb{K}$, the \\textit{active context} ${C_a}$ of $k_d$ is a subset ${C_a \\subseteq C}$ such that \n$\\forall k_i \\in C_a, \\nexists k_j \\in (C - C_a)~\\ni rel(k_j,k_d) > rel(k_i,k_d)$.\n\\end{definition}\nIn other words, $C_a$ contains the words in the input that are the most related ones to $k_d$. \nTo obtain such a context, we stick to the method proposed in~\\cite{gracia2009multiontology}:~1)~removing repeated words and stopwords from \\textit{C}, ~2) applying a semantic relatedness (${rel_{w}}$ in our case) between each context word ${k_i \\in C}$ and the keyword to disambiguate ${k_d}$, and ~3) constructing ${C_a}$ with the context words whose relatedness scores above a certain threshold. The output of this process is the \\textit{active context} ${C_a \\subseteq C}$. The maximum cardinality of ${C_a}$ is set to a fixed value(${|C_a| = 4}$) following Kaplan's experiments~\\cite{kaplan1955experiment}.\n\nOnce we have obtained $C_a$ for $k_d$, we\ncan apply the main algorithm, which takes as input $k_d$, $C_a$, and a set of possible senses for $k_d$, ${S_{k_d}}$.\nThe main steps are presented in Algorithm~\\ref{alg:disambiguation}\\footnote{We refer the interested reader to~\\cite{gracia2009multiontology} for the complete details.}: \n\n\n\\begin{algorithm}\n\n\\caption{Keyword disambiguation algorithm}\\label{alg:disambiguation}\n \\SetKwInOut{Input}{Input}\n \\SetKwInOut{keyword}{$K_d$}\n \\SetKwInOut{senses}{$S_{k_d}$}\n \\SetKwInOut{context}{$C_a$}\n \\SetKwInOut{Output}{Output}\n \n \\Input{}\n \\keyword{The keyword to disambiguate.}\n \\senses{The set of possible senses for $K_d$.}\n \\context{The active context selected for $K_d$.}\n \\Output{A weight for each sense ${s_i \\in S_{k_d}}$.}\n \\textbf{function} disambiguate $(K_d, S_{k_d}, C_a)$:\\\\\n \\ForEach{sense $s_i \\in S_{k_d}$}{\\label{line:semrelfrom}\n \\ForEach{keyword $k_j \\in C_a$}{\n $r_j=rel(s_i,k_j)$\n }\n $score_{s_i}=\\sum_{j}r_j\/|C_a|$\n }\\label{line:semrelto}\n $maxScore=max(score_{s_1},\\ldots,score_{s_n})$ \\label{line:overlapfrom}\\\\\n \\ForEach{sense $s_i \\in S_{k_d}$}{\n \n $newScore=score_{s_i}+(1-maxScore)*overlap(C_a,OC(s_i))$\n $score_{s_i}=newScore$\\label{line:overlap}\n }\\label{line:overlapto}\n $maxScore=max(score_{s_1},\\ldots,score_{s_n})$\\label{line:freqfrom} \\\\\n \\ForEach{sense $s_i \\in S_{k_d}$}{\n \\If{$score_{s_i} > proximityFactor*maxScore$}{\n $newScore=score_{s_i}+(1-maxScore)*normFreq(s_i)$\n $score_{s_i}=newScore$\n }\n }\\label{line:freqto}\n\\end{algorithm}\n\n\\begin{enumerate}\n \\item \\textbf{Applying the semantic relatedness:} First, the algorithm computes an initial disambiguation between the senses in ${S_{k_d}}$ and the active context $C_a$ (Lines~\\ref{line:semrelfrom}-\\ref{line:semrelto}). For this, we use the updated relatedness measure\n presented in the previous section (Equations~\\ref{eq:termWord} and~\\ref{eq:arccos}). The score assigned to each sense (${score_{s_i}}$) is the mean of $rel(s_i,k_j)$ where ${s_i \\in S_{k_d}}$ is a candidate sense of the keyword being disambiguated, and ${k_j \\in C_a}$ is a keyword in the active context.\n \\item \\textbf{Calculating the context overlap:} The disambiguation algorithm weights the scores taking into account the overlap between $C_a$ and the\n ontological context of each sense, ${OC(s_i)}$ (Lines~\\ref{line:overlapfrom}-\\ref{line:overlapto}).\n Note that ${OC(s_i)}$ includes its synonyms, glosses, and labels, as well as labels of other related terms, such as hypernyms, hyponyms, meronyms, holonyms, etc. The overlap is calculated (ignoring stopwords) as:\n $$overlap(C_a,OC(s_i)) =\\frac{|OC(s_i)\\cap C_a|}{min(|{OC(s_i)}|,|C_a|)}$$\n \n \n \n \n \\item \\textbf{Frequency of usage:} Finally,\n the frequency of use of the highest scored senses is taken into account (Lines~\\ref{line:freqfrom}-~\\ref{line:freqto}), if such information is available. The proximity decision is handled by a ${proximityFactor\\in[0,1]}$, which is combined with the maximum of the scores of the senses ($proximityFactor * maxScore$) to obtain a threshold. The scores of the senses $s_i$ which are above such a threshold are then updated using:\n \n \n $$normFreq({s_i})=\\sqrt{a*\\frac{frequency_{s_i}}{\\sum_{j}frequency_{s_j}}+b}$$ \n \n where $\\sum_{j}frequency_{s_j}$ is equal to the sum of the frequency of all senses of $k_d$, $a$ and $b$ are constrained\\footnote{We set $a = b = 0.5$ and $proximityFactor = 0.75$ as indicated in~\\cite{gracia2009multiontology}.} by $a,b \\in [0,1]$ and $a + b = 1$.\n\\end{enumerate}\n\nThe output of the disambiguation algorithm is a\nscore for each possible sense ${s_i \\in S_{k_d}}$ that represents the confidence level of being the right sense according to the active context $C_a$.\nNote that, in our approach, ${S_{k_d}}$ is not restricted to any particular dictionary, as it could be dynamically built from, e.g., different ontological resources.\n\n\\subsection{Proposed Modifications}\n\\label{sec:proposedmodifs}\nAs our aim is to study the best way to exploit word embeddings, we have analyzed their characteristics and explored different approaches to use them in the adopted disambiguation process. \nIn particular, in this section, we present a list of possible modifications to the Step~2 of the algorithm (Lines~\\ref{line:overlapfrom}-\\ref{line:overlapto}) to include and take advantage of the properties of word embeddings\nalong with the rationale behind them. For the rest of the section, let ${maxScore}$ be the maximum score among all senses in ${S_{k_d}}$, ${centroid}$ a function to calculate an average vector by the arithmetic mean of a set of vectors, and ${rel_{w}}$ the angular distance in Equation~\\ref{eq:arccos}. \nThus, the different approaches are described below:\n\n\\begin{itemize}\n \\item \\textbf{Average:} The straightforward way to include the embeddings is to calculate directly the average vector of the different bag of words involved in the disambiguation, under the assumption that the semantically coherent groups of words should outstand from the others. Thus, instead of computing the overlap between the semantic descriptions $OC(s_i)$ of each sense $s_i$ and the current active context $C_a$, we propose to compute the average between the word vectors from ${C_a}$ and ${OC(s_i)}$ to obtain a new score. Line~\\ref{line:overlap} in Algorithm~\\ref{alg:disambiguation} changes to:\n $$\\hspace{0.80cm}newScore(s_i)=score_{s_i}+(1-maxScore)*average({C_a},{OC(s_i)})$$\n where ${average({C_a},{OC(s_i)})}$ is:\n \\begin{align*}\n \\begin{split}\n average({C_a},{OC(s_i)})=\\frac{\\sum_{i, j}rel_{w}(k_i,w_j)}{|{C_a}||{OC(s_i)}|}, \\\\\n (k_i \\in {C_a}, w_j \\in {OC(s_i)}, i=1..|{C_a}|, j=1..|{OC(s_i)}|)\n \\end{split}\n \\end{align*}\n That is, we consider each set of words as a cluster in the vector space, and we represent them by their centroid. If there are elements that do not contribute to the semantic cohesion of the clusters, they will contribute negatively (they will increase the semantic distance) to select a particular sense for the target keyword\\footnote{We also studied other cluster-based distances measures (e.g., single linkage), but the results did not improve using the centroid-based measure, so we focused on the average vector which is broadly used in the literature.}.\n \n \\item \\textbf{Sense centroid without most frequent component:} As an evolution of the previous method, we studied the method described by Arora et al.~\\cite{arora2016simple}, called \\textit{Smooth Inverse Frequency (SIF)}. They propose to represent a sentence by a weighted average vector of its word vectors which the most frequent component using PCA\/SVD is substracted from. Thus, we propose to consider the semantic description ${OC(s_i)}$ of all senses of the sense inventory as sentences, and to calculate the SIF embedding of them. Then, during the disambiguation, we compute a new score (Line~\\ref{line:overlap}) for the sense being considered ${s_i}$ by measuring the distance between the centroid of the active context ${C_a}$ and the SIF vector of each $OC(s_i)$, following this computation:\\vspace{0.10cm}\n \n $\\hspace{0.25cm}newScore(s_i)=score_{si}+(1-maxScore)*$\n \n $\\vspace{0.15cm}\\hspace{2.30cm}rel_{w}(centroid({C_a}), SIF({OC(s_i)}))$\n \n Note that we do not substract the SIF vector from $C_a$ as all its words are already deemed as important.\n The most frequent component vector we are removing may encompass those words that occur most frequently in a corpus and lack semantic content (e.g., stop-words),\n thus not contributing to the actual disambiguation.\n \n \\item \\textbf{Top-K nearest words:} As a variant of the two previous methods, in this method, we select the top \\emph{k} nearest words from the semantic description ${OC(s_i)}$ of a sense to ${C_a} \\cup k_d$. After that, we compute the distance between centroids of the active context and the top K nearest words selected to obtain its new score: \\vspace{0.1cm}\n \n $\\hspace{0.25cm}newScore(s_i)=score_{s_i}+(1-maxScore) *$\n \n $\\hspace{2.30cm}rel_{w}(centroid({C_a}),centroid($\n \n $\\vspace{0.15cm}\\hspace{2.8cm}topKNearest(centroid({C_a}),{OC(s_i)})))$\n \n \n In this case, we work under the same hypothesis as for the selecing an \\emph{active context}: the words that belong to the semantic description of the sense that are the closest ones to the active context and the keyword that is being disambiguated,\n should be the most significant to contribute in making a correct disambiguation. \n \n \\item \\textbf{Doc2vec:} Finally, instead of treating the ontological descriptions as bag of words in this method, we consider them as proper documents and apply \\textit{doc2vec}~\\cite{le2014distributed}. In particular, each semantic description ${OC(s_i)}$ of the senses becomes a document, and \\textit{doc2vec} allows to calculate an embedding space for all of them. Then, we compute the distance between the centroid of the active context ${C_a}$ and the embedding calculated for the sense. Note that \\textit{doc2vec} learns as well a word embeddings model that it uses during training. We use those word vectors to create the centroid of the active context. Therefore, in a similar way, the new score is computed as:\\vspace{0.10cm}\n \n $\\hspace{0.25cm}newScore(s_i)=score_{s_i}+(1-maxScore)*$\n \n $\\vspace{0.15cm}\\hspace{2.30cm}rel_{w}(centroid({C_a}), doc2vec({OC(s_i)}))$\n \n We consider the semantic descriptions as documents\n to capture the distributional semantics both at local (window) and global (document) scope.\n \n\\end{itemize}\n\nWe report the best results that we obtained by applying these different approaches in the following section.\n\n\\section{Experimental Evaluation}\n\\label{sec:experimentaleval}\nIn this section, we discuss the results obtained in the experiments that we have carried out to evaluate our proposal.\nFirstly, we evaluated different available embedding models using the distance proposed in Equation~\\ref{eq:arccos}. We performed several tests comparing to human judgment in order to check how the \\emph{angular distance} behaved. Secondly, we evaluated the potential of our keyword disambiguation algorithm and the relatedness measure among ontology terms and words in the context of WSD, including all algorithm variations that we have proposed in Section~\\ref{sec:proposedmodifs}.\n\nFor the experiments, we used the following pre-trained vectors: \\textit{word2vec} trained on Google News corpus\\footnote{\\url{https:\/\/code.google.com\/archive\/p\/word2vec\/}}, \n\\textit{word2vec} trained on Wikipedia\\footnote{\\label{laugithub}\\url{https:\/\/github.com\/jhlau\/doc2vec}}, \\textit{doc2vec} trained also on Wikipedia\\footnote{Dump dated in 2015-12-01.}, \\textit{GloVe} trained on Wikipedia 2014 and Gigaword 5 corpus\\footnote{\\url{https:\/\/nlp.stanford.edu\/projects\/glove\/}}, and ${NASARI_{embed}}$ along with the \\textit{word2vec} word embeddings trained on the UMBC corpus\\footnote{\\url{http:\/\/lcl.uniroma1.it\/nasari\/\\#two}}. We used WordNet\\textsuperscript{\\ref{wordnetfootnote}} as the sense inventory.\n\n\\subsection{Correlation with Human Judgment}\n\\label{sec:correlation}\nIn order to validate the hypothesis of the suitability of using word embeddings along with the use of the \\emph{angular distance} to compute semantic relatedness, we first analysed the correlation of such a technique with human judgment in a basic word-to-word comparison. For this purpose, we used\ndifferent datasets available in the literature which contain pairs of words whose relatedness was manually assessed by different people. The datasets and their details can be seen in Table~\\ref{tab:datasets}.\n\n\n\\begin{table}[htb!]\n \\caption{Correlation with human judgment benchmarks.}\n \\centering\n \\begin{tabular}{|c|c|c|c|}\n \\hline\n Dataset & \\#Word Pairs & \\#Human Judges \\\\\n \\hline\n MC-30~\\cite{miller1991contextual} & 30 & 38 \\\\\n WordS353-Rel~\\cite{finkelstein2001placing} & 252 & 13 \\\\\n WordS353-Sim~\\cite{finkelstein2001placing} & 203 & 16 \\\\\n RG-65~\\cite{rubenstein1965contextual} & 65 & 51 \\\\\n MEN dataset (train\/dev)~\\cite{bruni2014multimodal} & 2000\/1000 & crowdsourced\\footnotemark \\\\\n GM dataset~\\cite{wise08} & 30 & 30 \\\\\n \\hline \n \\end{tabular}\n \\label{tab:datasets}\n\\end{table}\n\\footnotetext{They used Amazon Mechanical Turk: \\url{https:\/\/www.mturk.com\/}}\n\nThe results obtained for Spearman correlation are presented in Table~\\ref{tab:spearmancorrel}. Reported values, where available, were calculated using the widespread used cosine similarity. We can see that using the angular distance (Equation~\\ref{eq:arccos}) to calculate relatedness between pairs of words also offers a semantic correlation with the human judgment.\nIn particular, regarding the GM dataset~\\cite{wise08}, the authors reported a 78\\% using the previous \\emph{relWeb} measure (Equation~\\ref{eq:relweb}). We can see a strong improvement in this dataset by using word embeddings: we achieve up to a 87.3\\% using \\textit{word2vec} trained on Google News (taking into account the average of all models, we achieve an average of 81.2\\% for this dataset). These results enable us to use the angular distance as the core relatedness measure in Equations~\\ref{eq:completeeq} to~\\ref{eq:level1}. Note that for \\textit{word2vec} and \\textit{doc2vec} trained on Wikipedia we can not provide a comparison, because Lau \\& Baldwin~\\cite{lau-baldwin-2016-empirical} did not evaluate the correlation with human judgment.\n\n\\begin{table*}[!htb]\n \\caption{Spearman correlation coefficients between the angular distance applied on word pairs and human judgment in different datasets. Upper values are our evaluations, lower ones are the reported values in the original papers using the cosine distance. Highlighted values equal or outperform the best result (78\\%) for the same dataset in~\\cite{wise08}.}\n \\label{tab:spearmancorrel}\n \\begin{tabular}{|c|c|c|c|c|c|c||c|}\n \\hline\n Vectors\\textbackslash Datasets & MC-30& WS353-Sim& WS353-Rel& RG-65& MEN& GM & Average\\\\\n \\hline\n GloVe & 70.4& 66.5& 56.1& 76.9& 74.2& \\textbf{84.5} & 71.4\\\\\n \\hdashline\n Reported at~Pennington et al.~\\cite{pennington2014glove} & 72.7 & 65.8 & - & 77.8 & - & - & 72.1\\\\\n \\hline\n Google News word2vec& 80.0 & 77.2& 63.5& 76.0& 77.0& \\textbf{87.3} & 76.8\\\\\n \\hdashline\n Reported at Camacho-Collados et al.~\\cite{camacho2016nasari} & 80.0 & 77.0 & - & - & - & - & 78.5\\\\\n \\hline\n ${NASARI_{embed}+UMBCw2v}$& 70.3& 72.7& 56.8& 70.7& 74.5& 74.7 & 70.0\\\\\n \\hdashline\n Reported at Camacho-Collados et al.~\\cite{camacho2016nasari} & 83.0 & 68.0 & - & 80.0 & - & - & 75.5\\\\\n \\hline\n Wikipedia word2vec& 80.9& 77.9& 62.2& 78.3& 76.9& \\textbf{81.8} & 76.3\\\\\n \\hdashline\n Reported at Lau \\& Baldwin~\\cite{lau-baldwin-2016-empirical} & - & - & - & - & - & - & -\\\\\n \\hline\n Wikipedia doc2vec& 73.3& 69.0& 52.3& 71.6& 72.0& \\textbf{77.8} & 69.3\\\\\n \\hdashline\n Reported at Lau \\& Baldwin~\\cite{lau-baldwin-2016-empirical} & - & - & - & - & - & - & -\\\\\n \\hline\n \\end{tabular}\n\\end{table*}\n\n\\subsection{Word Sense Disambiguation Evaluation}\nTo evaluate our proposal, we used three datasets oriented to WSD: SemCor 2.0 dataset\\footnote{\\url{http:\/\/web.eecs.umich.edu\/\\~mihalcea\/downloads.html\\#semcor}}, SemEval2013 all-words WSD dataset\\footnote{\\label{semeval2013footnote}\\url{https:\/\/www.cs.york.ac.uk\/semeval-2013\/task12.html}}, and SemEval2015 All-Words Sense Disambiguation and Entity Linking dataset\\textsuperscript{\\ref{semeval2015footnote}}.\nWe used WordNet\\textsuperscript{\\ref{wordnetfootnote}} as sense inventory. For SemCor~2.0, we specifically used WordNet~2.0 as such a dataset is annotated with this version, and for the rest of datasets we used WordNet 3.0.\n\nWe tested all the options proposed in Section~\\ref{sec:proposedmodifs} for the disambiguation algorithm, and we obtained that the \\textit{Top-K nearest words} option achieves the best results\\footnote{Achieving the best performance for ${K=15}$.}. Thus, due to space restrictions, we focus on \\textit{Top-K nearest words} option in this section\\footnote{\\label{expsurl}The interested reader can find all the details of the experiments at \\url{https:\/\/bit.ly\/2lqCzop}}. \nRegarding the models, we selected \\textit{word2vec} trained in Google News and \\textit{word2vec} trained in Wikipedia \nbecause they showed better average correlation with human judgment in different datasets (see Table~\\ref{tab:spearmancorrel}); and ${NASARI_{embed}+UMBCw2v}$ word embeddings because, although they do not excelled in correlation with human judgment, they showed the best performance in all test datasets for WSD. Finally,\nin order to compare the results to~\\cite{wise08}, we report the precision results for SemCor 2.0; while we report the F-score results for the rest of datasets.\n\n\\subsubsection*{SemCor2.0 Experiments:}\nFollowing~\\cite{wise08}, in this set of experiments, for each of three selected highly ambiguous nouns (\\textit{plant}, \\textit{glass}, and \\textit{earth}), we took~10 random sentences from the corpus. Table~\\ref{tab:semcor2} presents the results: all cases outperform the results achieved in~\\cite{wise08}, which reported an average precision of 57\\%. Our best performance is an average precision of 63.15\\% with ${NASARI_{embed}+UMBCw2v}$ vectors.\nIn fact, SIF method shows equal or even slightly better performance in this particular dataset using ${NASARI_{embed}+UMBCw2v}$ vectors and \\textit{word2vec} trained in Google News vectors. However, in the rest of cases it is \\textit{Top-K nearest words} method that obtains the best results. \nIn addition, SIF method requires to preprocess the target sense inventory to calculate the sentence embeddings, introducing a mild dependence to it.\nOur selected method shows a good performance (it improves the results of the original algorithm), while allowing to be more decoupled from the actual sense inventory used.\n\n\\begin{table*}[!htb] \n \\caption{Precision results for SemCor 2.0 dataset (10 random sentences) adopting Top-K nearest words. The two rightmost columns show the results using the \\textit{relWeb} based relatedness measure and Most Frequent Sense methods as reported in~\\cite{wise08}.}\n \\label{tab:semcor2}\n \\begin{tabular}{|c|c|c|c||c|c|}\n \\hline\n Experiment\\textbackslash Approach& Wikipedia word2vec & Google News word2vec & ${NASARI_{embed}+UMBCw2v}$ & \\textit{relWeb}* & Most Freq. Sense*\\\\\n \\hline\n 10 sent. with PLANT& 58.44\\%& 63.03\\%& 66.20\\%& 80\\% & 40\\%\\\\\n 10 sent. with GLASS& 57.47\\%& 63.78\\%& 60.15\\%& 30\\% & 30\\%\\\\\n 10 sent. with EARTH& 59.21\\%& 56.38\\%& 62.33\\%& 60\\% & 60\\%\\\\\n \\hline\n AVERAGE& 58.41\\%& 61.13\\%& 63.15\\%& 57\\%& 43\\%\\\\ \n \\hline\n \\end{tabular}\n\\end{table*}\n\n\\subsubsection*{SemEval Results:} In Table~\\ref{tab:fscore}, we present the results obtained for SemEval 2013 and SemEval 2015. In this case, ${NASARI_{embed}+UMBC}\\\\{w2v}$ vectors achieved the best results, with an F-score of 64.39\\%. In SemEval 2013, UMCC-DLSI reported the best results, with a F-score of 64.7\\%, similar to ours. Besides, our results are similar to other state-of-the-art systems using sense embeddings\\KILL{vectors}: Camacho et al.~\\cite{camacho2016nasari} reported an F-Score value of 66.7\\% with their ${NASARI_{lexical}}$ vectors evaluated in SemEval 2013. Unfortunately, they do not provide results for their ${NASARI_{embed}}$ vectors using WordNet.\nRegarding SemEval 2015, the best reported result in our task reached an F-score of 65.8\\%, while the baseline used to compare systems (BabelNet First Sense (BFS)) was an F-score of 67.5\\%. We reach an F-score value of 61.61\\%, which, while does not beat previous values, is meritory given that our approach is focused on situations where linguistic information might be scarce (e.g., keyword-based input).\n\n\\begin{table*}[!htb]\n \\caption{F-Score results for SemEval 2013 and 2015 datasets adopting Top-K nearest words\\KILL{and using WordNet 3.0 as sense inventory}. Values in column highlighted by ${\\diamond}$ correspond to the best system in each SemEval dataset. Values in column highlighted by ${^\\pm}$ are the baselines reported in SemEval 2013~and SemEval 2015.}\n \\label{tab:fscore}\n \\begin{tabular}{|c|c|c|c||c|c|}\n \\hline\n & Wikipedia word2vec & Google News word2vec & ${NASARI_{embed}+UMBCw2v}$ & Best system${^\\diamond}$ & Baseline${^\\pm}$ \\\\\n \\hline\n SemEval 2013 & 59.59\\% & 62.81\\% & 64.39\\% & 66.7\\% & 63.0\\%\\\\\n SemEval 2015 & 61.32\\% & 60.37\\% & 61.61\\% & 65.8\\% & 67.5\\%\\\\\n \\hline\n \\end{tabular}\n\\end{table*}\n\nTo sum up, our proposal improves the results presented in~\\cite{wise08} by substituting their Web search engine-based measure to one that uses word embeddings. We also improve the disambiguation results reported in~\\cite{gracia2009multiontology} by adapting their algorithm to exploit the properties of the word embeddings. Our proposal achieves similar performance levels to the SOTA, while providing the flexibility to work independently of the resources used (i.e., word embeddings, sense inventory), and reducing the barriers to its application to any domain.\n\n\\section{Conclusions and future work}\n\\label{sec:conclusionsandfuturework}\nIn this paper, we have presented a keyword disambiguation approach based on a semantic relatedness measure which exploits the semantic information captured by word embeddings and tries to map the meanings from a sense inventory. We have visited the semantic relatedness measure proposed in~\\cite{wise08} to adapt it to work with word embeddings instead of relying on Web search engines, and we have improved a disambiguation algorithm~\\cite{gracia2009multiontology} by exploring different uses and types of embeddings (both at word and sentence level).\n\nTo validate our proposal, we have performed several experiments around Word Sense Disambiguation (WSD) where we have used different pre-trained word embeddings and WordNet as the resource to obtain the target senses of words. With our proposal:\n\\begin{itemize}\n \\item We are able to relate words and meanings from a sense inventory (e.g., ontology terms) in a flexible way, by exploiting available resources and regardless the domain in which we are working. This makes our measure adaptive and general enough to be used for different contexts. \n \\item We provide a method which can be adapted to any domain in a dictionary-decoupled way, provided that we have a document corpus which would allow us to capture the distributional semantics. This lowers the requirements of data in order to build more specific models for particular domains. \n \\item We have tested the capabilities of different word embedding models, improving the results presented in~\\cite{wise08}. We evaluated our measure in the same SemCor 2.0 dataset described in this work, and we have obtained in the best case an average increase of 6\\% in precision (a relative improvement of about~11\\%).\n \\item Being decoupled from a fixed pool of senses does not come at the expense of performance. We achieve similar quality of results than having an ad hoc and more expensive trained model capturing the possible senses. In particular, we have tested our measure in SemEval 2013 and SemEval 2015 datasets reaching an F-score of 64.39\\% and 61.61\\% respectively. These results are similar to the state of the art~\\cite{camacho2016nasari} using \\textit{sense2vec} approaches.\n\\end{itemize}\n\nAs future work, we want to further extend our approach to the field of concept discovery (similar to entity search~\\cite{Balog:2018:ER}, but focused on concepts rather than on instances). We also want to explore newer contextualized word embeddings, such as BERT or XLNet (based in ELMo~\\cite{peters2018deep}), and how they could be used in this context.\nFinally, we would like to propose a specific dataset exclusively for keyword disambiguation taking QALD\\footnote{QALD is a series of evaluation campaigns on Question Answering over Linked Data: (\\url{http:\/\/qald.aksw.org}).} datasets as baseline; we want to develop it in order to test our relatedness measure in a more appropriate dataset for the context in which we focus: the disambiguation of keyword-based inputs.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n Finding a new chaotic system began from the last six decades to understand different natural and man made systems. Still scientific communities are investigating new chaotic system due to its practical applications in engineering \\cite{szemplinska1991}, finance \\cite{guegan2009}, plasma \\cite{ray2008, ray2009}, time series analysis \\cite{kodba2005, nazarimehr2018a} and natural observations like, climatology \\cite{selvam2007}, biology \\cite{may1987, ghosh2007}, geology \\cite{turcotte1997} etc. After discovery of Lorenz model \\cite{lorenz1963} in 1963, there are several chaotic models were investigated like, R\\\"{o}ssler system \\cite{rossler1976}, Chen system \\cite{chen1999}, Chua's attractor \\cite{chua1986}, Sprott system \\cite{sprott1994, sprott2000}, logistic model \\cite{may1976}, predator-prey based model \\cite{hastings1991, ginoux2019} and many more. At first, theoretical background behind generation of chaotic solution in deterministic system was not developed. But, now different ways are brought to light to understand and realize the chaotic attractors and its dynamics, both from analytical as well as numerical approaches \\cite{eckmann1985, perc2005, silva2018}. Initially, three well-known routes were reported for generating chaos in dynamical systems, $\\it i.e.$, period-doubling, quasi-periodic and intermittent routes \\cite{ott2002}. Nowadays, researchers are fascinated for finding chaotic attractor in a different manner. One interesting topic is to find multistable chaotic attractors \\cite{bayani2019} which has important application and impact to explore a broad field of research \\cite{dudkowski2016}. Besides a new concept of attractor arises connected with multistability like attractors with no equilibria \\cite{wei2011, jafari2013b, pham2014}, line equilibria \\cite{jafari2013a, nazarimehr2018b}, plane of equilibrium \\cite{jafari2016} or coexistence of more than one chaotic attractors \\cite{dudkowski2016, ray2009} etc. Basin of attraction of chaotic attractor does not touch with small neighborhood of the unstable fixed point, which is quite different from the traditional chaotic attractors (called {\\it self-excited chaotic attractors}). This kind of attractor is familiar as {\\it hidden attractor} as finding of hidden attractor is difficult task usually \\cite{dudkowski2016}. Researches are also eager to control of the multistability \\cite{pisarchik2014}. But finding a simple chaotic system is still a challenging task with a great impact for basic research \\cite{liu2004, lu2002, wang2012}. We have also knowledge about chaotic systems in which different timescales are separated which have huge impact in various fields, like chemical oscillatory model \\cite{epstein1983}, Hindmarsh-Rose neuronal model \\cite{hindmarsh1984}, prey-predator model \\cite {kuwamura2009} and so on. Well known features of chaotic attractor are sensitive dependence of initial conditions, boundedness, aperiodic and appearance of it is confirmed by existence of positive Lyapunov exponent in the parameter region \\cite{strogatz2014}. Researchers are also very much interested to study the chaotic dynamics in delay systems due to its impactful applicability \\cite{wei2019, khajanchi2018, banerjee2019, ghosh2017}. With studying about generation of chaotic system, parallel researchers have been focused about control of chaos in non-delay chaotic systems \\cite{ott1990, perc2006, jun2008} as well as delay chaotic systems \\cite{ghosh2008, bhowmick2014}. Its application includes chemical reaction oscillations, turbulent fluids etc. \\cite{boccaletti2000}. From the above discussion, the finding of a new chaotic system and exploring different dynamical states deserve a special attention.\n\n\\par In this paper, we propose a simple nonlinear system which exhibits single-scroll chaotic attractor in a certain parameter region. At first, we calculate the fixed points of the system and analyze their linear stability analyses. From these analyses, we observe that the two fixed points out of three always unstable, whereas one fixed point changes its stability via Hopf bifurcation and finally goes to chaotic state through period-doubling route. We also analytically derive the first Lyapunov coefficient to check the stability of the limit cycle. Then, we introduce slow-fast ratio parameter in the proposed model and investigate the previously mentioned system parameter region by varying this slow-fast parameter. We analytically calculate the slow manifold with the help of geometric singular perturbation theory. The chaotic behavior completely annihilate if the time scale ratio parameter is low.\n\\par The rest of the paper is organized as follows: Section-2 devotes the introduction of a new chaotic system and linear stability analysis of the fixed points. We also discuss the bifurcation diagram, Lyapunov exponent and two-dimensional parameter regions by varying the system parameters. In Sec-3, we discuss the effect of slow-fast parameter and controlling of chaos. Finally, conclusions are drawn in Sec-4. \n\n\\section{Model and system dynamics}\nWe propose a three-dimensional model with two nonlinear turns in the following form,\n\\begin{equation}\n\\begin{array}{l}\\label{eq.1}\t\n\\dot{x} =-x+y+z,\\\\\n\\dot{y} =xy-z,\\\\\n\\dot{z} =-axz+y+b,\\\\\n\\end{array}\n\\end{equation}\nwhere $a,~b$ are two positive parameters and $x, y$ and $z$ are state variables. Here dot denotes time derivative with respect to time $t$. We first study the linear stability analysis of the equilibrium points of the above system. The equilibrium points are in the form $(x^*, f(x^*), g(x^*))$, where $f(x)$ and $g(x)$ are given by\n\\[\nf(x)=\\frac{x}{1+x},~~~g(x)=\\frac{x^2}{1+x}.\n\\]\nHere $x^*$ satisfies the following equation\n\\begin{equation}\n\\begin{array}{l}\\label{eq.2}\t\nax^*{^3}-(b+1)x^*-b=0.\n\\end{array}\n\\end{equation}\nSolving Eqn. (\\ref{eq.2}), we have three values of $x^*$ corresponding to three roots $x_1^*$, $x_2^*$ and $x_3^*$, \n\\begin{equation*}\n\\begin{array}{l}\nx_1^*=2\\sqrt{\\frac{b+1}{3a}}cos(\\frac{\\theta}{3}),\\\\\nx_2^*=-2\\sqrt{\\frac{b+1}{3a}}cos(\\frac{\\pi-\\theta}{3}),\\\\\nx_3^*=-2\\sqrt{\\frac{b+1}{3a}}cos(\\frac{\\pi+\\theta}{3}),\n\\end{array}\n\\end{equation*}\n\\begin{figure}[ht]\n\t\\centerline{\n\t\t\\includegraphics[scale=0.6]{figure_1.png} \t}\n\t\\caption{(a) $a$-$b$ parameter space where region A denotes the existence of three fixed points, whereas one fixed point exists in the region B. (b) Variation of maximum real part of the eigenvalues corresponding to the three fix points by varying the parameter $a$. Here maximum eigenvalues corresponding to two points $x_1^*$ and $x_2^*$ are always positive (red and magenta curves) and $x_3^*$ is negative up to $a<1.53$ and change the sign after that which is shown in inset figure. Here one fixed point changes its stability and other two are always unstable. Here $b=0.25$.} \n\t\\label{fig1} \n\\end{figure}\nwhere\n\\[\n\\theta=\\tan^{-1}\\Big(\\frac{1}{3b}\\sqrt{\\dfrac{4(b+1)^3-27ab^2}{3a}}\\Big).\n\\]\nThe above three roots of the Eqn. (\\ref{eq.2}) are real if \n\\[\n00$, Eqn.\\ (\\ref{eq.3}) gives one real root with a pair of complex conjugate eigenvalues and another case, $i.e.$, $p^3+q^2<0$, the equation gives three real roots.\n\n\\par Now, we fix the system parameter $b$ at $b=0.25$ and consider $a$ as a bifurcation parameter. In Fig.\\ \\ref{fig1}(b), we plot the maximum real part of the eigenvalues of $J$ at $x_1^*, x_2^*,$ and $x_3^*$ with respect to the system parameter $a$. From this figure, the maximum real part of the eigenvalue (Re(${\\lambda_{3}}$)) of $J$ at third root, $i.e.,~x_3^*$ shows negative in the range $a \\in$ (0, 1.53) (blue line) and after that turns to positive (clearly shown in the inset figure of Fig.\\ \\ref{fig1}(b)), but for remaining two roots, $i.e.,~x_1^*,~x_2^*$, their associated maximum eigenvalues (Re(${\\lambda_{1,2}}$)) are always positive (magenta and red lines). Interestingly, we see that at $a=4.63$, there are one maximum real part of the eigenvalues corresponding to $x_1^*, x_2^*,$ and $x_3^*$ exists. This means for $a>4.63$, the Eqn. (\\ref{eq.2}) possesses only one real root ($B$ region in Fig.\\ \\ref{fig1}(a)). So it can be concluded that the Eqn. (\\ref{eq.1}) has only one stable fixed point corresponding to $x_3^*$ in the range $a \\in (0, 1.53)$. Now, for a specific value of $a=0.5$ and $b=0.25$, we calculate $x^*$ and associated eigenvalues $\\lambda_{1, 2, 3}$ as follows:\n\\begin{equation*}\n\t\\begin{array}{l}\n\t\tx_1^*: x^*\\approx1.6731,\\\\ (\\lambda_{1},\\lambda_{2},\\lambda_{3})\\approx(1.6791,-0.9213\\pm0.9524i);\\\\\n\t\tx_2^*: x^*\\approx-1.4698,\\\\ (\\lambda_{1},\\lambda_{2},\\lambda_{3})\\approx(-3.0764,1.718,-0.3765);\\\\ \n\t\tx_3^*: x^*\\approx-0.2034,\\\\ (\\lambda_{1},\\lambda_{2},\\lambda_{3})\\approx(-0.9653,-0.0682\\pm1.1072i). \n\t\\end{array}\n\\end{equation*}\n\n\\begin{figure*}[ht]\n\t\\centerline{\n\t\t\\includegraphics[scale=0.31]{figure_2_a.png}\n\t\t\\includegraphics[scale=0.31]{figure_2_b.png}\n\t\t\\includegraphics[scale=0.31]{figure_2_c.png}}\n\t\\caption{(a) Bifurcation diagram of the new system Eqn.\\ (\\ref{eq.1}) for $a\\in[2.4,3.037]$ showing period-doubling route to chaos.\n\t\t(b)\tZoom version of a periodic window shown as shaded region in figure (a). Here LPC, PD, and IC denote the limit point bifurcation of cycles, period-doubling, and interior crisis, respectively. (c) Corresponding maximum Lyapunov exponent with respect to $a\\in[2.4,3.037]$ to verify different states (periodic state and chaotic state). Here, we fix $b=0.25$. Initial condition is taken as $(x(0), y(0), z(0))=(0.3, 0.2, 0.1).$}\n\t\\label{fig2} \n\\end{figure*}\n\\begin{figure}[ht]\n\t\\centerline{\n\t\t\\includegraphics[scale=0.55]{figure_3_r.png}}\n\t\\caption{Phase portraits for different values of $a$: $a=2.4$, period-1 (magenta); $a=2.6$, period-2 (black); $a=2.75$, period-4 (blue) and $a=2.9$, chaotic state (red). From the phase diagrams, period-doubling route to chaos is easily observed. Here $b=0.25$. }\n\t\\label{fig3} \n\\end{figure}\nUsing the sign of eigenvalues, the fixed points corresponding to $x_1^*, x_2^*,$ and $x_3^*$ of the system (\\ref{eq.1}) are unstable saddle focus, unstable saddle and stable focus, respectively at $a=0.5$ and $b=0.25$. So we focus on the fixed point corresponding to $x_3^*$ and study Hopf bifurcation analysis. The characteristic equation at the point $x_3^*$ becomes,\n\\begin{equation}\n\\begin{array}{l}\\label{eq.5}\n\\lambda^3+r_2\\lambda^2+r_1\\lambda+r_0=0,\n\\end{array}\n\\end{equation}\nwhere\n\\begin{equation*}\n\t\\begin{array}{l}\n\t\tr_2=1+(1-a)(\\frac{\\sqrt{15a}}{3a})cos(\\theta_1),\\\\\n\t\tr_1=\\dfrac{1-\\frac{(a-1)\\sqrt{15a}}{3a}cos(\\theta_1)+\\frac{5(a-1)}{3a}cos^2\\theta_1+\\frac{5\\sqrt{15a}}{9a}cos^3\\theta_1}{1-(\\frac{\\sqrt{15a}}{3a})cos(\\theta_1)},\\\\\n\t\tr_0=\\dfrac{1-5cos^2(\\theta_1)+\\frac{10\\sqrt{15a}}{9a}cos^3(\\theta_1)}{1-(\\frac{\\sqrt{15a}}{3a})cos(\\theta_1)},\n\t\\end{array}\n\\end{equation*}\nand,\n\\[\n\\theta_1=\\frac{\\sqrt{3a(125-27a)}}{9a}.\n\\]\nAccording to the Routh-Hurwitz stability criterion, the real parts of all the eigenvalues $\\lambda$ become negative if\n$r_2>0, r_0 >0,$ and $r_1r_2 >r_0$. Numerically, we calculate the Hopf bifurcation point at $a^* \\approx 1.524053$ where the real parts of the complex conjugate roots of Eqn. (\\ref{eq.5}) becomes zero \\cite{kuznetsov2004}.\n\\begin{figure}[ht]\n\t\\centerline{\n\t\t\\includegraphics[scale=0.65]{figure_4_r.png} \t}\n\t\\caption{ Basin of attraction of chaotic attractor at (a) $z=0.1$ plane, (b) $y=0.1$ plane, and (c) $x=0.1$ plane. White region represents the unbounded state. (d) First return map for the chaotic attractor. This map shows unimodal behavior which signify period-doubling route to chaos. Initial condition is chosen at $x(0)=0.3, y(0)=0.2,$ and $z(0)=0.1$. Parameter values are fixed at $a=2.9$ and $b=0.25$}\n\t\\label{fig4} \n\\end{figure} \n\nNow, we discuss about the nature of Hopf bifurcation \\cite{kuznetsov2004}. For this purpose, first we check the transversality condition and then calculate the value of first Lyapunov coefficient \\cite{kuznetsov2004, mello2009, wei2014}. The transversality condition, $\\ i.e.$, real part of $\\frac{d\\lambda(a^*)}{da}|_{\\lambda=1.0843i} \\approx 0.06839 (\\ne 0)$ is satisfied here. The first Lyapunov coefficient is $l_1\\approx0.04196~(> 0)$ (detail calculations are in appendix). So the fixed point (corresponding to $x_3^*$) deals with subcritical Hopf bifurcation at $a^* \\approx 1.524053$ and it indicates that there is an unstable limit cycle exists in the neighborhood of the fixed point.\n\n\\par Next, we study the dynamical behavior of the equilibrium point corresponding to $x_3^*$ with respect to the parameter $a$ after crossing the value $a=a^*$. After crossing $a^*$, the fixed point loses its stability and a stable limit cycle of period-1 appears and persists up to $a\\approx2.493$. Then we achieve chaotic dynamics via period-doubling (PD) route \\cite{strogatz2014}. This scenario is depicted in Fig.\\ \\ref{fig2}(a). At $a\\approx3.037$, this chaotic attractor is destructed due to boundary crisis \\cite{grebogi1983} as the chaotic attractor collides with unstable fixed point (corresponding to $x_1^*$). We take a zoom version of periodic window at the neighborhood of $a=2.845$ (shaded region in Fig.\\ \\ref{fig2}(a)) in Fig.\\ \\ref{fig2}(b). System changes from chaotic oscillation to periodic oscillation through limit point bifurcation of cycles (LPC). Chaotic oscillation again revives via PD bifurcation and amplifies via interior crisis (IC). Then we plot maximum Lyapunov exponent ($\\rm LE_{max}$) in Fig.\\ \\ref{fig2}(c) for verifying our result in Fig.\\ \\ref{fig2}(a). Positive value of $\\rm LE_{max}$ after a certain value of $a$ signifies the presence of chaos. For four different values of $a$, (i.e., $a=2.4, 2.6, 2.75$ and $2.9$), we display four different kind of attractors, $\\ i.e.$, limit cycle of period-1 (magenta), period-2 (black), period-4 (blue), and chaos (red), respectively. The occurrence of period-doubling cascade for generating chaos is verified by calculating Feigenbaum universal constant \\cite{strogatz2014}. The ratio of distances between successive intervals between the bifurcate values of the system parameter, represented by a constant factor\n\\[\n\\delta=\\lim_{k\\to\\infty} \\frac{a_{k+1}-a_{k}}{a_{k+2}-a_{k+1}}\n\\] \napproaches to 4.66992.... The value of Feigenbaum constant for our case are given in the following table: \n\\vspace{0.1cm}\n\\begin{table}[h!]\n\t\\centering\n\t\\begin{tabular}{||c c c c||} \n\t\t\\hline\n\t\tk & Period & Bifurcation parameter $a_k$ & Ratio $\\frac{a_{k+1}-a_{k}}{a_{k+2}-a_{k+1}}$ \\\\ [1ex] \n\t\t\\hline\\hline\n\t\t1 & $2^1$ & 2.491 & Not applicable \\\\ \n\t\t2 & $2^2$ & 2.7234 & Not applicable \\\\\n\t\t3 & $2^3$ & 2.7822 & 3.952 \\\\\n\t\t4 & $2^4$ & 2.7966 & 4.0833 \\\\\n\t\t5 & $2^5$ & 2.79969 & 4.6602 \\\\\n\t\t6 & $2^6$ & 2.8003517 & 4.6696 \\\\ [1ex] \n\t\t\\hline\n\t\\end{tabular}\n\\end{table}\\vspace{1.0cm}\n\\begin{figure}[ht]\n\t\\centerline{\n\t\t\\includegraphics[scale=0.45]{figure_5.png} \t}\n\t\\caption{ Different dynamical behaviors in the $a$-$b$ parameter space. Here Cyan: stable steady state, magenta: period-1, black: period-2, blue: period-4, and red: greater than period-4 or chaotic state. White region is for unbounded solution. The initial condition is $x(0)=0.3, y(0)=0.2$, and $z(0)=0.1$.}\n\t\\label{fig5} \n\\end{figure}\n\n Next, we identify the region of initial conditions for the chaotic state. The basin of attractions at $z=0.1$ plane, $y=0.1$ plane, and $x=0.1$ plane for the chaotic attractor at $(a, b)=(2.9, 0.25)$ are depicted in Figs.\\ \\ref{fig4}(a-c). We get fully connected basin of attractions for respective three cases. At the same set of parameter value, we also plot first return map, $\\ i.e.$, $x_n$ vs. $x_{n+1}$ in Fig.\\ \\ref{fig4}(d), where $x_{n}$ denotes the $n$-th local maxima of the time series. Unimodal nature of this first return map indicates the PD bifurcation route to generate chaos. \n\\par To explore the complete study by varying the two parameters $a$ and $b$ simultaneously, we plot a two parameter $(a$-$b)$ bifurcation diagram \\cite{ghosh2015} for the range $a \\in [0,5]$ and $b \\in [0.01,0.5]$ in Fig.\\ \\ref{fig5}. This demonstrates the whole dynamics of the system (\\ref{eq.1}). In this phase diagram, cyan, magenta, black, blue, and red regions represent steady state, limit cycle of period-1, period-2, period-4, and greater than period-4 (or chaotic), respectively. White portion represents unbounded solution for higher values of $a$ and $b$. Interestingly, well-separated parameter space is observed in $(a$-$b)$ bifurcation diagram. \n \n\\section{Slow-fast dynamical system and controlling chaos}\n In this section, we introduce a slow-fast time scale ratio parameter to create slow-fast dynamics in the proposed system. We try to understand the effect of this time scale parameter. A general framework of a slow-fast dynamical system \\cite{guckenheimer2008} can be written as,\n\t\\begin{equation*}\n\t\\begin{array}{lll}\n\t\t\t\t\n\t\\epsilon \\dfrac{dx}{dt}=f(x,y,\\eta); \n\t~\\dfrac{dy}{dt}= g(x,y,\\eta),\n\t\\end{array}\n\t\\end{equation*}\n\twhere $x\\in \\mathbb{R}^{m} \\;\\mbox{and}\\; y\\in \\mathbb{R}^{n}$ are fast and slow variables, respectively, $\\eta\\in \\mathbb{R}^{p}$ is a model parameter, $0<\\epsilon<1$ represents the ratio of time-scales. \tSetting $\\epsilon$=0, the trajectory of above equation converges to the solution of differential algebraic equation $f(x,y,\\eta) =0$ and $\\dot {y}=g(x,y,\\eta)$, where \t$S_{0}= \\{(x,y) \\in \\mathbb{R}^{m} \\times \\mathbb{R}^{n}| f(x,y,\\eta) =0\\}$ is a critical manifold.\n\n\n\\par Now, in our model we introduce a time-scale ratio parameter $\\epsilon ~(0<\\epsilon<1)$ in the Eqn. (\\ref{eq.1}) as follows \n\\begin{equation}\n\\begin{array}{l}\\label{eq.6}\t\n\\epsilon \\dot{x} =-x+y+z,\\\\\n\\dot{y} =xy-z,\\\\\n\\dot{z} =-axz+y+b,\\\\\n\\end{array}\n\\end{equation}\nso that $x$ is fast variable and $y, z$ are slow variables. We rescale the above Eqn. (\\ref{eq.6}) by setting $\\tau=t\/\\epsilon$ as \n\\begin{equation}\n\\begin{array}{l}\\label{eq.7}\t\nx^{'} =-x+y+z,\\\\\ny^{'} =\\epsilon(xy-z),\\\\\nz^{'} =\\epsilon(-axz+y+b),\\\\\n\\end{array}\n\\end{equation}\nwhere ${x^{'}}\\equiv\\frac{dx}{d\\tau}$. By setting $\\epsilon=0$, the trajectory of Eqn. (\\ref{eq.6}) converges to the solution of the differential algebraic equation \n\\begin{equation}\n\\begin{array}{l}\\label{eq.8}\t\n0 =-x+y+z,\\\\\n\\dot{y} =xy-z,\\\\\n\\dot{z} =-axz+y+b.\\\\\n\\end{array}\n\\end{equation}\n\n\\begin{figure}[ht]\n\t\\centerline{\n\t\t\\includegraphics[scale=0.55]{figure_7.png} }\n\t\\caption{Plot of the critical manifold $S_0$ and slow manifold $S_{\\epsilon}$ at $\\epsilon=0.85$ are illustrated by green and red colors, respectively. Chaotic attractor (blue line) is drawn for $a=3.17$ and $b=0.26$. The position of the unstable fixed point (solid dot): $(-0.2420, -0.2286, 0.0553)$.}\n\t\\label{fig6} \n\\end{figure}\n\nSo, as per above definition, $S_0=\\{(x,y,z)\\in \\mathbb R^3 : x=y+z\\}$ is called a critical manifold. From Fenichel theory \\cite{fenichel1979}, we know that a subset $S_{\\epsilon}$ of $S_{0}$ exists for $\\epsilon~(>0)$ if $S_{\\epsilon}$ is normally hyperbolic. Our system is normally hyperbolic because of real part of $\\frac{\\partial (-x+y+z)}{\\partial x}=-1 (\\ne 0), (x,y,z)\\in S_{\\epsilon}$. So there exists a local invariant manifold (slow manifold) $S_{\\epsilon}$ and we derive this equation by using geometric singular perturbation method \\cite{fenichel1979, guckenheimer2008, desroches2012}. \n\\par In $S_0$, relation among variables $x,y$ and $z$ is $x=y+z$. Let in $S_{\\epsilon}$, this relation becomes\n\\begin{equation}\n\\begin{array}{l}\\label{eq.9}\t\nx=y+z+O(\\epsilon).\n\\end{array}\n\\end{equation}\nBy Taylor series expansion in $\\epsilon$, we get \n\\begin{equation}\n\\begin{array}{l}\\label{eq.10}\t\nx=y+z+\\epsilon H(y,z)+O(\\epsilon^2).\n\\end{array}\n\\end{equation}\nDifferentiate Eqn. (\\ref{eq.9}) and then multiply both sides by $\\epsilon$, we have,\n\\begin{equation}\n\\begin{array}{l}\\label{eq.11}\n\\epsilon \\dot{x}=\\epsilon \\dot{y}+\\epsilon \\dot{z}+O(\\epsilon^2).\n\\end{array}\n\\end{equation}\nUsing Eqns. (\\ref{eq.6}) and (\\ref{eq.9}) in Eqn. (\\ref{eq.11}), we get \n\\begin{equation}\n\\begin{array}{l}\\label{eq.12}\n\\epsilon \\dot{x}=\\epsilon\\Big (y^2-az^2+(1-a)yz+y-z+b\\Big )+O(\\epsilon^2).\n\\end{array}\n\\end{equation}\nAgain from first equation of Eqn. (\\ref{eq.6}), \n\\[\n\\epsilon \\dot{x} =-x+y+z.\n\\]\nNow using Eqn. (\\ref{eq.10}) in above equation, we get \n\\begin{equation}\n\\begin{array}{l}\\label{eq.13}\n\\epsilon \\dot{x}=-(\\epsilon H(y, z)+O(\\epsilon^2)).\n\\end{array}\n\\end{equation}\nComparing Eqns. (\\ref{eq.12}) and (\\ref{eq.13}), finally we obtain\\\\\n\\begin{equation}\n\\begin{array}{l}\\label{eq.14}\nH(y,z)=-\\Big (y^2-az^2+(1-a)yz+y-z+b\\Big ).\n\\end{array}\n\\end{equation}\n\\begin{figure*}[ht]\n\t\\centerline{\\includegraphics[scale=0.05]{figure_8_a_1.png}\n\t\t\\includegraphics[scale=0.5]{figure_8_a_2.png}}\n\t\\centerline{\\includegraphics[scale=0.6]{figure_8_b.png}}\n\t\\caption{Upper panel: different dynamical states in $a$-$b$ parameter spaces for (a) $\\epsilon=0.9$, and (b) $\\epsilon=0.1$. Here, Cyan: stable steady state, magenta: period-1, black: period-2, blue: period-4, and red: greater than period-4 or chaotic state. White region is for unbounded solution. Lower panel: bifurcation diagrams with respect to the time-scale ratio parameter $\\epsilon$ for (c) $\\{a, b\\}=\\{3.17, 0.26\\}$ and (d) $\\{a, b\\}=\\{2.68, 0.304\\}$. The initial condition is $x(0)=0.3, y(0)=0.2$, and $z(0)=0.1$. From this figure, it is noticed that chaotic state is annihilated by decreasing the value of time-scale parameter $\\epsilon$ and stable steady state and periodic solutions dominate the whole parameter region.}\n\t\\label{fig7} \n\\end{figure*} \n\\par We use the above value of $H(y, z)$ in Eqn. (\\ref{eq.10}) to calculate the approximated slow manifold $S_{\\epsilon}$. We plot the critical manifold $S_0$ (in green color), an approximated slow manifold $S_{\\epsilon}$ (in red color) in Fig.\\ \\ref{fig6}. The chaotic attractor for the set of values $a=3.17, b=0.26$ and $\\epsilon=0.85$ is shown by blue line. This attractor lies on the both sides of the manifold $S_0$ and $S_{\\epsilon}$. We take the value $\\epsilon=0.85$ for $S_{\\epsilon}$ (though this value is not so small enough, but we choose this value due to the appearance of chaos). \n\\par Next we study the effect of time-scale parameter $\\epsilon$ of the system (\\ref{eq.6}) in ($a$-$b$) parameter space for the previously mentioned range. For this, we first set $\\epsilon=0.9$ and different dynamical states like, stable fixed point, period-1, period-2, period-4 and greater than period-4 or chaos are shown in Fig.\\ \\ref{fig7}(a). The stable fixed point dominates in the ($a$-$b$) parameter space compare to Fig.\\ \\ref{fig5} for $\\epsilon=1.0$ case, i.e., without any time-scale. Also the region of chaotic state decreases compare to the system (\\ref{eq.1}) (in Fig.\\ \\ref{fig5}). If we decrease the value of time-scale parameter $\\epsilon$, then the chaotic region is completely annihilated in the ($a$-$b$) parameter space. Figure \\ref{fig7}(b) shows that the region of chaotic state is completely annihilated whereas region of stable fixed point is expanded and rest parts of the region are filled by period-1 limit cycle for $\\epsilon=0.1$. So by introducing the slow-fast parameter $\\epsilon$ in our original system (\\ref{eq.1}), we can control chaos. Now we take two sets of parameter values of $\\{a, b\\}, \\ i.e.,$ $\\{3.17, 0.26\\}$ and $\\{2.68, 0.304\\}$ from the Figs.\\ \\ref{fig7}(a, b), we have illustrated two bifurcation diagrams by varying slow-fast ratio parameter $\\epsilon \\in[0,1]$ for understanding the various change of dynamics in slow-fast system in Figs.\\ \\ref{fig7}(c, d), respectively. For the first set of parameter values, the period-doubling route to chaos is observed, whereas for second case, the chaotic behavior in the system is terminated and two states (stable fixed point and period-1 limit cycle) are dominated. \n\\section{Conclusions}\n\\par To summarize, in this paper, we have proposed a new simple chaotic system with two nonlinear terms. We have elaborately discussed the nature of fixed points through linear stability analysis. We have also analytically derived the first Lyapunov coefficient to show the nature of Hopf bifurcation by which fixed point changes its stability. Using bifurcation diagram, we have investigated the period-doubling route to chaos and also verified the results using Lyapunov exponent and first return map. We also separate the region of different dynamical behavior in the parameter space. Finally, we introduce slow-fastness in our model and try to understand the regarding changes of dynamics. Variation of dynamics has been captured in ($a$-$b$) parameter space for different values of the slow-fast parameter $\\epsilon$. Most interesting fact is that we can control the chaotic and higher periodic orbits by adjusting the value of $\\epsilon$\n\\\\\n\\par {\\bf Acknowledgments:}\nThe authors would like to thank Syamal K. Dana, Chittaranjan Hens, Gourab Kumar Sar, Sayeed Anwar, Soumen Majhi, Srilena Kundu, Sayantan Nag Chowdhury, Subrata Ghosh and Sarbendu Rakshit for helpful discussions and comments. \n\\section*{Appendix: Calculation of first Lyapunov coefficient} \n\\par Let us consider the differential equation \n\\begin{equation}\n\\begin{array}{l}\\label{eq.15}\n\\dot{x}=f(x,\\alpha),\n\\end{array}\n\\end{equation}\nwhere $x\\in \\mathbb{R}^3$ and $\\alpha\\in \\mathbb{R}^1$ represent the state variables and system parameter respectively, and $f:\\mathbb{R}^3\\times \\mathbb{R}^1 \\rightarrow \\mathbb{R}^1$ is a smooth function. Let the Eqn. (\\ref{eq.15}) has a fixed point $x=x_0$ at $\\alpha=\\alpha_0$. Now, replacing $x-x_0$ by $X$ and Eqn. (\\ref{eq.15}) becomes\n\\begin{equation}\n\\begin{array}{l}\\label{eq.16}\n\\dot{X}=F(X)=f(X,\\alpha_0).\n\\end{array}\n\\end{equation}\nHere, $F(X)$ is also a smooth function and we can expand it in Taylor series in terms of symmetric functions as follows,\n\\begin{equation}\n\\begin{array}{l}\\label{eq.17}\nF(X)=AX+\\frac{1}{2!}B(X,X)+\\frac{1}{3!}C(X,X,X)+O(||X||^4),\n\\end{array}\n\\end{equation} \nwhere $A=f_x(0,\\alpha_0)$ is the Jacobian matrix and for i=1,2,3, \n\\begin{equation}\n\\begin{array}{l}\\label{eq.18}\nB_{i}(X,Y)=\\sum_{j,k=1}^{3} \\frac{\\partial^2{F_i(\\xi)}}{\\partial \\xi_{j} \\partial\\xi_{k}}|_{\\xi=0} X_j X_k,\\\\\nC_{i}(X,Y,Z)=\\sum_{j,k,l=1}^{3} \\frac{\\partial^3{F_i(\\xi)}}{\\partial \\xi_{j} \\partial\\xi_{k} \\partial\\xi_{l}}|_{\\xi=0} X_j X_k X_l.\n\\end{array}\n\\end{equation} \nIn our case, multilinear symmetric functions are\\\\\n\\begin{equation}\n\\begin{array}{l}\\label{eq.19}\nB(X,Y)=(0,x_1y_2+x_2y_1,-a^*(x_1y_3+x_3y_1)),\\\\\nC(X,Y,Z)=(0,0,0).\n\\end{array}\n\\end{equation}\n\n\\par We assume that A has a pair of complex eigenvalues $\\lambda_{2,3}=\\pm iw_0(w_0>0)$ on the imaginary axis at $(x_0,\\alpha_0)$. Let generalized eigen space of $A$ be $T^C$ which is the largest invariant subspace spanned by eigenvectors corresponding to $\\lambda_{2,3}$. Let $p,q\\in$ C$^3$ be such eigen vectors then \n\\begin{equation}\n\\begin{array}{l}\\label{eq.20}\nAq=iw_0q,\\\\\nA^Tp=iw_0p,\\\\\n=\\sum_{i=1}^{3}\\bar{p}_iq_i=1,\n\\end{array}\n\\end{equation}\nwhere $A^T$ denotes the transpose of the matrix $A$ and $<...>$ is the standard scalar product over $\\mathbb{C}^3$. For our system (\\ref{eq.1}), we obtain\\\\\n$\nq\\approx(-0.0404 -0.5130i,~ 0.6402,~ -0.1244 - 0.5568i),\\\\\np \\approx(-0.1629 - 0.0591i,~ 0.78 - 0.1677i,~ 0.1501 - 0.8436i),\\\\\n$\nwhere $w_0\\approx 1.0843$.\n\\par Now we choose any $y\\in T^C$ which is represented as $y=wq+\\bar w\\bar q$, where $w=\\in \\mathbb{C}$. The two-dimensional center manifold associated with the eigenvalues $\\lambda_{2,3}= \\pm iw_{0}$ can be parameterized by $w$ and $\\bar w$ in the form $X=H(w,\\bar w)$. Here $H: \\mathbb{C}^2\\rightarrow\\mathbb{R}^3$ can be expanded as follows\n\\begin{equation}\n\\begin{array}{l}\\label{eq.21}\nH(w,\\bar w)=wq+\\bar w\\bar q\t+\\sum_{2\\le j+k\\le 3}\\frac{1}{j!k!}h_{jk}w^j\\bar w^k+O(|w|^4),\n\\end{array}\n\\end{equation}\nwhere $h_{jk}\\in \\mathbb{C}^3$ and $h_{jk}=\\bar h_{kj}$. Differentiating $H(w,\\bar w)$ with respect to $t$ and substituting this into Eqn. (\\ref{eq.15}) and using Eqn. (\\ref{eq.16}), we get the following differential equation\n\\begin{equation}\n\\begin{array}{l}\\label{eq.22}\nH_{w}\\dot{w}+H_{\\bar w}\\dot{\\bar w}=F(H(w,\\bar w)).\n\\end{array}\n\\end{equation} \nUsing Eqns. (\\ref{eq.16}), (\\ref{eq.21}) and (\\ref{eq.22}), we determine the complex coefficients $h_{ij}$ and from the Eqn. (\\ref{eq.22}), the $w$ evolves on the center manifold,\\\\\n\\begin{equation}\n\\begin{array}{l}\\label{eq.23}\n\\dot{w}=iw_0w+\\frac{1}{2}G_{21}w|w|^2+O(|w|^4),\n\\end{array}\n\\end{equation}\nwhere $G_{21}\\in \\mathbb{C}$.\nSubstituting Eqn. (\\ref{eq.23}) into Eqn. (\\ref{eq.22}) and using Eqns.\\~ (\\ref{eq.17}) and (\\ref{eq.18}), we get the following coefficients,\n\\begin{equation}\n\\begin{array}{l}\\label{eq.24}\nh_{11}=-A^{-1}B(q,\\bar q),\\\\\nh_{20}=(2iw_0I_3-A)^{-1}B(q,q),\n\\end{array}\n\\end{equation}\nwhere $I_3$ is the 3 $\\times$ 3 identity matrix. Finally $G_{21}$ is obtained from the condition of existence of solution of the equation for $h_{21}$. So, \n\\begin{equation*}\n\t\\begin{array}{l}\n\t\tG_{21}=,\n\t\\end{array}\n\\end{equation*}\nwhere we use $=1$. So the first Lyapunov coefficient is defined as \n\\[\nl_1=\\frac{1}{2}Re(G_{21}),\n\\]\nwhich determines the nonlinear stability of a nondegenerate codimension one Hopf bifurcation.\nNow, for calculating $l_1$, we first find \\\\\n$ h11\\approx( 0.6348, 1.0854, -0.4507),$\\\\\n$h20\\approx( -0.2273 + 0.0747i, -0.1643 - 0.1240i, -0.2251 - 0.2942i).$\nusing Eqns. (\\ref{eq.18}), (\\ref{eq.24}) and (\\ref{eq.24}).\nFinally, we get $l_1\\approx0.04196~(> 0)$. \n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzbwde b/data_all_eng_slimpj/shuffled/split2/finalzzbwde new file mode 100644 index 0000000000000000000000000000000000000000..0c821adbc91dc0ebcacda87b59b565ceee92e4b3 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzbwde @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\nThe European Union (EU) is now composed from 27 countries \nand is considered as a major world leading power \\cite{eusite}.\nJanuary 2021 has seen Brexit officially taking place, triggering the \nwithdrawal of the United Kingdom (UK) from EU \\cite{wikibrexit}.\nThis event has important political, economical and social effects.\nHere we project and study its consequences from the view point\nof international trade between world countries.\nOur analysis is based on the UN COMTRADE database \\cite{comtrade} for\nthe multiproduct trade between world countries in \nrecent years. From this database we construct the world trade\nnetwork (WTN) and evaluate the influence and trade power \nof specific countries using the Google matrix analysis\nof the WTN. We consider 27 EU countries as a single\ntrade player having the trade exchange between EU and other countries.\nOur approach uses the Google matrix tools and algorithms developed for the WTN \n\\cite{wtn1,wtn2,wtn3,keu9} and other complex directed networks\n\\cite{rmp2015,politwiki}. The efficiency of the Google matrix and PageRank algorithms\nis well known from the World Wide Web network analysis \\cite{brin,meyer}.\n\nOur study shows that the Google matrix approach (GMA) allows to characterize \nin a more profound manner the trade power of countries \ncompared to the usual method relying on import and export analysis (IEA) between countries.\nGMA's deeper analysis power originates in the fact that it accounts for the multiplicity\nof transactions between countries while IEA only takes into account the \neffect of one step (direct link or relation) transactions. In this paper, we show that the \nworld trade network analysis with GMA identifies EU as the first trade player in the world, \nwell ahead of USA and China.\n\nThis paper is structured in the following way. Section~\\ref{sec:dataset} introduces first the UN COMTRADE dataset, \nand then gives a primer on the tools related to Google matrix analysis such as the trade balance metric and the REGOMAX algorithm.\nIn Section~\\ref{sec:results}, the central results of this papers are presented, which are discussed in \\ref{sec:discussion}. \n\n\\section{Data sets, algorithms and methods}\\label{sec:dataset}\n\nWe use the UN COMTRADE data \\cite{comtrade} for years 2012, 2014, 2016 and 2018 to \nconstruct the trade flows of the multiproduct WTN following the procedure\ndetailed in \\cite{wtn2,wtn3}. This paper gives the results for year 2018 only, others are \nto be found at \\cite{ourwebpage}. Each year is presented by a money matrix, \n$M^{p}_{cc^{\\prime}}$, giving the export flow of product $p$\nfrom country $c^{\\prime}$ to country $c$ (transactions are expressed in USD of current year).\nThe data set is given by $N_{c} = 168$ countries and territories \n(27 EU countries are considered as one country) and $N_{p} = 10$ principal type of products\n(see the lists in \\cite{wtn1,wtn3}). These 10 products are: Food and live animals (0);\nBeverages and tobacco (1); Crude materials, inedible, except fuels (2); Mineral fuels etc (3);\nAnimal and vegetable oils and fats (4); Chemicals and related products, n.e.s. (5);\nBasic manufactures (6); Machinery, transport equipment (7); \nMiscellaneous manufactured articles (8);\nGoods not classified elsewhere (9) (product index $p$ is given in brackets). They belong to \nthe Standard International Trade Classification (SITC Rev. 1)\nThus the total Google matrix $G$ size is given by all system nodes $N=N_c N_p= 1680$\nincluding countries and products.\n\nThe Google matrix $G_{ij}$ of direct trade flows\nis constructed in a standard way described in detail at \\cite{wtn2,wtn3}: monetary trade flows \nfrom a node $j$ to node $i$ are normalized to unity \nfor each column $j$ thus given the matrix $S$ of \nMarkov transitions for trade, the columns of dangling nodes with\nzero transactions are replaced by a column with all elements being $1\/N$. \nThe weight of each product is taken into account via a\ncertain personalized vector taking into account \nthe weight of each product in the global trade volume.\nWe use the damping factor $\\alpha=0.5$. \nThe Google matrix is $G_{ij}=\\alpha S_{ij} + (1-\\alpha) v_i$\nwhere $v_i$ are components of positive column vectors called personalization vectors\nwhich take into account the weight of each product in the global trade \n($\\sum_i v_i=1$).\nWe also construct the matrix $G^*$ for the inverted trade flows.\n\nThe stationary probability distribution \ndescribed by $G$ is given by the PageRank vector\n$P$ with maximal eigenvalue $\\lambda=1$:\n$GP=\\lambda P =P$ \\cite{rmp2015,brin,meyer}. \nIn a similar way, for the inverted flow, described by $G^*$, \nwe have the CheiRank vector $P^*$, \nbeing the eigenvector of $G^* P^* = P^*$. \n PageRank $K$ and CheiRank $K^*$ indexes are obtained from\nmonotonic ordering of probabilities of PageRank vector $P$ and \nof CheiRank vector $P^*$ as\n$P(K)\\ge P(K+1)$ and $P^*(K^*)\\ge P^*(K^*+1)$ with $K,K^*=1,\\ldots,N$.\nThe sums over all products $p$ gives the PageRank and CheiRank \nprobabilities of a given country as $P_c =\\sum_p P_{cp}$ and \n${P^*}_c =\\sum_p {P^*}_{cp}$ (and in a similar way\nproduct probabilities $P_p, {P^*}_p$) \\cite{wtn2,wtn3}.\nThus with these probabilities we obtain the related indexes $K_c, {K^*}_c$.\nWe also define \nfrom import and export trade volume the\n probabilities\n$\\hat{P}_p$, $\\hat{P}^*_p$, $\\hat{P}_c$, $\\hat{P}^*_c$,\n$\\hat{P}_{pc}$, $\\hat{P}^*_{pc}$ and \ncorresponding indexes\n$\\hat{K}_p$, $\\hat{K}^*_p$, $\\hat{K}_c$, $\\hat{K}^*_c$, $\\hat{K}$, $\\hat{K}^*$\n(these import and export probabilities\nare normalized to unity\nby the total import and export volumes, see details in\n\\cite{wtn2,wtn3}). It is useful to note that qualitatively\nPageRank probability is proportional to the volume of ingoing\ntrade flow and CheiRank respectively to outgoing flow.\nThus, we can approximately consider that the\nhigh import gives a high PageRank $P$ probability\nand a high export a high CheiRank $P^*$ probability.\n\n\\begin{figure}[!ht]\n\t\\centering\n\t\\includegraphics[width=0.99\\textwidth]{fig1.pdf}\n\t\\caption{Circles with country flags show country positions on\n the plane of PageRank-CheiRank indexes $(K,K^*)$\n (summation is done over all products) (left panel)\n and on the plane of ImportRank-ExportRank $\\hat{K}$, $\\hat{K}^*$ \n from trade volume (right panel);\n data is shown only for index values less than $61$ in year 2018\n \n .} \\label{fig1}\n\\end{figure}\n\n\nAs in \\cite{wtn2,wtn3},\nwe define the trade balance of a given country\nwith PageRank and CheiRank probabilities given by\n$B_c = (P^*_c - P_c)\/(P^*_c + P_c)$.\nAlso we have from ImportRank and ExportRank probabilities\nas $\\hat{B}_c= ({\\hat{P}^*}_c - \\hat{P}_c)\/({\\hat{P}^*}_c + \\hat{P}_c)$.\nThe sensitivity of \ntrade balance $B_c$ to the price of energy or machinery can be obtained \nfrom the change of corresponding money volume\nflow related to SITC Rev.1 code $p=3$ (mineral fuels) or $p=7$ (machinery)\nby multiplying it by $(1+\\delta)$, renormalizing column to unity and \ncomputing all rank probabilities and\nthe derivatives $dB_c\/d\\delta$. \n\nWe also use the REGOMAX algorithm \\cite{politwiki,wtn3}\nto construct the reduced Google matrix $G_R$ \nfor a selected subset of WTN nodes $N_r \\ll N$.\nThis algorithm takes into accounts all transitions of direct and\nindirect pathways happening in the full Google matrix $G$ \nbetween $N_r$ nodes of interest. We use this $G_R$ matrix\nto construct a reduced network of most strong transitions\n(``network of friends'')\nbetween a selection of nodes representing countries and products.\n\nEven if Brexit enter into play in 2021, we use UN COMTRADE data of previous years\nto make a projecting analysis of present and future power of EU composed of 27 countries.\n\nFinally we note that GMA allows to obtain interesting results for \nvarious types of directed networks including Wikipedia \n\\cite{geopolitics,celestin} and protein-protein interaction \n\\cite{zinprotein,signor} networks.\n\n\\begin{table}[]\n \\centering\n \\caption{Top 20 countries of PageRank ($K$), CheiRank ($K^*$), \n ImportRank and ExportRank in 2018.}\n \\begin{tabular}{|c|l|l|l|l|}\n \\hline\n Rank & PageRank & CheiRank & ImportRank & ExportRank \\\\ \\hline\n1 & EU & EU & EU & China \\\\\n2 & USA & China & USA & EU \\\\\n3 & China & USA & China & USA \\\\\n4 & United Kingdom & Japan & Japan & Japan \\\\\n5 & India & Repub Korea & United Kingdom & Repub Korea \\\\\n6 & U Arab Emirates & India & Repub Korea & Russia \\\\\n7 & Japan & Russia & India & United Kingdom \\\\\n8 & Mexico & U Arab Emirates & Canada & Mexico \\\\\n9 & Repub Korea & Singapore & Mexico & Canada \\\\\n10 & Canada & United Kingdom & Singapore & Singapore \\\\\n11 & Singapore & South Africa & Switzerland & Switzerland \\\\\n12 & Switzerland & Thailand & U Arab Emirates & India \\\\\n13 & Turkey & Malaysia & Russia & Malaysia \\\\\n14 & Russia & Canada & Thailand & Australia \\\\\n15 & Australia & Mexico & Viet Nam & U Arab Emirates \\\\\n16 & South Africa & Turkey & Australia & Thailand \\\\\n17 & Thailand & Australia & Turkey & Saudi Arabia \\\\\n18 & Brazil & Switzerland & Malaysia & Viet Nam \\\\\n19 & Saudi Arabia & Brazil & Indonesia & Brazil \\\\\n20 & Malaysia & Saudi Arabia & Brazil & Indonesia \\\\ \\hline\n \\end{tabular}\n\\label{tab1}\n\\end{table}\n\n\\section{Results}\\label{sec:results}\n\\subsection{CheiRank and PageRank of countries}\n\nThe positions of countries on the PageRank-CheiRank $(K,K^*)$ and\nImportRank-ExportRank $(\\hat{K},\\hat{K}^*)$ planes are shown in Fig.~\\ref{fig1}\nand in Table~\\ref{tab1}.\nThese results show a significant difference between these two types of ranking.\nIndeed, EU takes the top PageRank-CheiRank position $K=K^*=1$\nwhile with Export-Import Ranking it has \nonly $\\hat{K}=1; \\hat{K}^*=2$, with USA at $\\hat{K}=2, \\hat{K}^*=3$\nand China at $\\hat{K}=3, \\hat{K}^*=1$. Thus EU takes the leading positions\nin the GMA frame which takes into account the muliplicity of trade transactions\nand characterizes the robust features of EU trade relations.\nAlso GMA shows that UK position is significantly weakened compared to IEA description\n(thus UK moves from $K^*=7$ in IEA to $K^*=10$ in GMA). From this data, we see also examples\nof other countries that significantly improve there rank positions in GMA\nframe compared to IEA: India ($K=5$, $K^*=6$, $\\hat{K}=7$, $\\hat{K}^*=12$), \nUnited Arab Emirates ($K=6$, $K^*=8$, $\\hat{K}=12$, $\\hat{K}^*=15$), \nSouth Africa ($K=16$, $K^*=11$, $\\hat{K}=23$, $\\hat{K}^*=23$). \nWe attribute this to well developed, deep \nand broad trade network of these countries which are well captured by GMA in contrast\nto IEA. Indeed, IEA only measures the volume of direct trade exchanges, while GMA characterises \nthe multiplicity of trade chains in the world.\n\n\\begin{figure}[!h]\n\t\\centering\n\t\\includegraphics[width=\\textwidth]{fig2.pdf}\n\t\\caption{World map of trade balance of countries \n$B_c = ({P_c}^* - P_c)\/({P_c}^* + P_c)$.\nTop: trade balance values are computed from the trade volume of Export-Import;\nbottom: trade balance values are computed from PageRank\nand CheiRank vectors; $B_c$ values are marked by color with the corresponding \ncolor bar marked by $j$; \ncountries absent in the UN COMTRADE report are marked by black color \n(here and in other Figs).}\\label{fig2}\n\\end{figure}\n\n\\subsection{Trade balance and its sensitivity to product prices}\n\nThe trade balance of countries in IEA and GMA frames is shown in Fig.~\\ref{fig2}.\nThe countries with 3 strongest positive balance are: \nEquatorial Guinea ($B_c=0.732$), Congo ($B_c=0.645$), Turkmenistan ($B_c=0.623$) in IEA and\nChina ($B_c=0.307$), Japan ($B_c= 0.244$), Russia ($B_c= 0.188$) in GMA.\nWe see that IEA marks top countries which have no significant world power\nwhile GMA marks countries with real significant world influence.\nFor EU and UK we have respectively $B_c=-0.015; 0.020$ (EU) and \n$B_c=-0.178; -0.187$ (UK) in IEA; GMA.\nThus the UK trade balance is significantly reduced in GMA\ncorresponding to a loss of network trade influence of UK\nin agreement with data of Fig.~\\ref{fig1} and Table~\\ref{tab1}.\n(We note that the balance variation bounds in GMA are smaller compared to IEA;\nwe attribute this to the fact of multiplicity of transactions \nin GMA that smooth various fluctuations which are more typical for IEA).\n\n\\begin{figure}[t]\n\t\\includegraphics[width=\\textwidth]{fig3.pdf}\n\t\\caption{Sensitivity of country balance \n$dB_c\/d\\delta_s$ for product $s=3$ (mineral fuels).\nTop: probabilities are from the trade volume of Export-Import;\nbottom: probabilities are from PageRank\nand CheiRank vectors. Color bar marked by $j$ \ngives sensitivity.} \\label{fig3}\n\\end{figure}\n\nThe balance sensitivity $d B_c\/d \\delta_s$ to product $s=3$\n(mineral fuels (with strong petroleum and gas contribution))\nis shown in Fig.~\\ref{fig3}. The top 3 strongest \npositive sensitivities $d B_c\/d \\delta_s$ \nare found for Algeria (0.431), Brunei (0.415), South Sudan (0.411) in IEA\nand Saudi Arabia (0.174), Russia (0.161), Kazakhstan (0.126) in GMA.\nThe results of GMA are rather natural since Saudi Arabia, Russia and Kazakhstan are central petroleum producers. \nIt is worth noting that GMA ranks Iraq at the 4th position. \nThe 3 strongest negative sensitivities are Zimbabwe (-0.137), Nauru (-0.131), \nJapan (-0.106), in IEA\nand Japan (-0.066), Korea (-0.062), Zimbabwe (-0.058), in GMA.\nFor China, India we have $d B_c\/d \\delta_s$ values being respectively:\n-0.073, -0.086 in IEA and -0.056, 0.010 in GMA. This shows that the trade network\nof India is more stable to price variations of product $s=3$.\nThese results demonstrate that GMA selects more globally influential countries.\n\n\\begin{figure}[t]\n\t\\includegraphics[width=\\textwidth]{fig4.pdf}\n\t\\caption{Same as in Fig.~\\ref{fig3} but\nfor product $s=7$ (machinery).} \\label{fig4}\n\\end{figure}\n\nThe balance sensitivity $d B_c\/d \\delta_s$ to product $s=7$\n(machinery) is shown in Fig.~\\ref{fig4}. \nHere the top 3 strongest \npositive sensitivities $d B_c\/d \\delta_s$ are found\nin both IEA and GMA for Japan (respectively 0.167, 0.151), \nRepub. Korea (0.143, 0.097), Philippines (0.130, 0.091).\nThe 3 strongest negative sensitivities are Brunei (-0.210), Iran (-0.202), \nUzbekistan (-0.190) in IEA\nand Russia (-0.138), Kazakhstan (-0.102), Argentina (-0.097) in GMA.\nThus we again see that GMA selects more globally influential countries. \nThe sensitivity $d B_c\/d \\delta_s$ values for EU, UK, China, Russia, USA are:\n EU (0.048), UK (0.006), China (0.065), Russia (-0.170), USA (-0.027) in IEA;\n EU (0.000), UK (0.024), China (0.077), Russia (-0.138), USA (-0.018) in GMA.\nLatter GMA results show that even if machinery product ($s=7$) is very important for EU \nthe network power of trade with this product becomes dominated by Asian countries\nJapan, Repub. Korea, China, Philippines; \nin this aspect the position of UK is slightly better than EU. \n\n\\begin{figure}[t]\n\t\\includegraphics[width=\\textwidth]{fig5.pdf}\n\t\\caption{Sensitivity of country balance \n$dB_c\/d\\delta_{cs}$ for product price $s=7$ (machinery) \nof EU (top left), USA (top right), China (bottom left)\nand $s=3$ (mineral fuel) of Russia (bottom right);\n$B_c$ is computed from PageRank and CheiRank vectors;\nsensitivity values are marked by color with the corresponding \ncolor bar marked by $j$. \nFor EU, USA, China, Russia we have $dB_c\/d\\delta_{cs} = 0.11, 0.11, 0.14, 0.12$\nrespectively, these values are \nmarked by separate magenta color to highlight sensitivity \nof other countries in a better way.} \\label{fig5}\n\\end{figure}\n\n\\begin{figure}[t]\n\t\\includegraphics[width=\\textwidth]{fig6.pdf}\n\t\\caption{Network trade structure between \nEU, USA (US), China (CN), Russia (RU) \nwith 10 products. \nNetwork is obtained from the reduced\nGoogle matrix $G_R$ (left) and ${G^*}_R$ (right)\nby tracing four strongest outgoing links\n(similar to 4 ``best friends'').\nCountries are shown by circles with two letters of country\nand product index listed in Section 2. The arrow direction \nfrom node $A$ to node $B$ means that $B$ imports from $A$ (left)\nand $B$ exports to $A$ (right).\nAll $40$ nodes are shown.\n} \\label{fig6}\n\\end{figure}\n\nIn Figs.~\\ref{fig3} and \\ref{fig4}, we have considered the sensitivity of \ncountry balance to a global price of a specific product (mineral fuel $s=3$ or machinery $s=7$). \nIn contrast, with GMA, we can also obtain the sensitivity of country balance to the price of\nproducts originating from a specific country. Such results are shown in Fig.~\\ref{fig5}.\nThey show that machinery ($s=7$) of EU gives a significant positive balance sensitivity \nfor UK and negative for Russia. This indicates a strong dependence of Russia \nfrom EU machinery. Machinery of USA gives strong positive effect for Mexico and Canada\nwith a negative effect for EU, Russia, Brazil, Argentina.\nMachinery of China gives positive sensitivity for Asian countries\n(Repub. Korea, Japan, Philippines) and significant negative effect for Mexico.\nMineral fuels ($s=3$) of Russia gives positive effect for Kazakhstan, Uzbekistan, Ukraine\n(former USSR republics) and negative effect for competing petroleum and gas producers\nNorway and Algeria.\n\n\\subsection{Network structure of trade from reduced Google matrix}\n\nThe network structure for 40 nodes of 10 products of EU, USA, China and Russia is \nshown in Fig.~\\ref{fig6}. It is obtained from the reduced Google matrix of $N_r=40$ nodes\nof global WTN network with $N=1680$ nodes on the basis of REGOMAX algorithm\nwhich takes into account all pathways between $N_r$ nodes via the global network of $N$ nodes.\nThe networks are shown for the direct ($G$ matrix) and\ninverted ($G^*$ matrix) trade flows. For each node we show only 4 strongest outgoing \nlinks (trade matrix elements) that heuristically can be considered as the four ``best friends''. \nThe resulting network structure clearly shows\nthe central dominant role of machinery product. For ingoing flows (import direction) of $G_R$ \nthe central dominance of machinery for USA and EU is directly visible while \nfor outgoing flows (export direction), machinery of EU and China dominate exports.\n\nIt is interesting to note that the network influence of EU with 27 countries\nis somewhat similar to the one constituted by a kernel of 9 dominant EU countries (KEU9)\n(being Austria, Belgium, France, Germany, Italy,\nLuxembourg, Netherlands, Portugal, Spain) discussed in \\cite{keu9}. This shows\nthe leading role played by these KEU9 countries in the world trade influence of EU.\n\nFinally we note that additional data with figures and tables is available at \\cite{ourwebpage}.\n\n\\section{Discussion}\\label{sec:discussion}\n\nWe presented the Google matrix analysis of multiproduct WTN obtained from UN COMTRADE database\nin recent years. In contrast to the legacy Import-Export characterization of trade,\nthis new approach captures multiplicity of trade transactions between world countries\nand highlights in a better way the \nglobal significance and influence of trade relations between specific countries and products.\nThe Google matrix analysis clearly shows that the dominant position\nin WTN is taken by the EU of 27 countries despite the leave of UK after Brexit.\nThis result demonstrates the robust structure of worldwide EU trade.\nIt is in contrast with the usual Import-Export analysis in which USA and China\nare considered as main players. We also see that machinery and mineral fuels\nproducts play a dominant role in the international trade.\nThe Google matrix analysis stresses the growing dominance of \nmachinery products of Asian countries (China, Japan, Republic of Korea).\n\nWe hope that the further development of Google matrix analysis \nof world trade will bring new insights in this complex system of world economy.\n\n{\\it Acknowledgments:}\nWe thank Leonardo Ermann for useful discussions.\nThis research has been partially supported through the grant\nNANOX $N^\\circ$ ANR-17-EURE-0009 (project MTDINA) in the frame \nof the {\\it Programme des Investissements d'Avenir, France} and\nin part by APR 2019 call of University of Toulouse and by \nRegion Occitanie (project GoIA).\nWe thank UN COMTRADE for providing us a friendly access\nto their detailed database.\n\n\n\n\n\n\\section{#1}\n \\addcontentsline{lof}{section}{%\n Section \\thesection: #1 \\vspace{10pt}\n }\n}\n\n\n\n\n\\renewcommand{\\figurename}{SupFig.}\n\n\n\\date{February 2021}\n\n\\usepackage{graphicx}\n\\usepackage{hyperref}\n\n\\begin{document}\n\n\\title{Supplementary Material for: \\\\ Post-Brexit power of European Union \\\\ from the world trade network analysis}\n\n\\author{Justin Loye\\inst{1,2} \\and\nKatia Jaffr\\`es-Runser\\inst{3}\n\\and\nDima L. Shepelyansky\\inst{2}\n}\n\\authorrunning{J.Loye et al.}\n\\titlerunning{Supplementary Material for: Post-Brexit power of European Union from the WTN}\n\n\\institute{Institut de Recherche en Informatique de Toulouse, \\\\\nUniversit\\'e de Toulouse, UPS, 31062 Toulouse, France\\\\\n\\email{justin.loye@irit.fr}\\\\\n\\and\nLaboratoire de Physique Th\\'eorique, \nUniversit\\'e de Toulouse, CNRS, UPS, 31062 Toulouse, France\\\\\n\\email{dima@irsamc.ups-tlse.fr}\\\\\n\\and\nInstitut de Recherche en Informatique de Toulouse, \\\\\nUniversit\\'e de Toulouse, Toulouse INP, 31071 Toulouse, France\\\\\n\\email{katia.jaffresrunser@irit.fr}\\\\\n}\n\\maketitle \n\n\\vspace{1.0cm}\n\nWe present additional results SupFig.~\\ref{fig:KKstar_2012_2016}, \\ref{fig:balance2012_2016}, \\ref{fig:importations},\n\\ref{fig:exportations}, \\ref{fig:dbalance_petroleum_2012_2016} and \\ref{fig:dbalance_machinery_2012_2016}.\n\nSupFig.~\\ref{fig:KKstar_2012_2016}, \\ref{fig:balance2012_2016}, \\ref{fig:dbalance_petroleum_2012_2016} and \\ref{fig:dbalance_machinery_2012_2016}\nreproduce for years 2012, 2014, 2016 respectively Fig.1, Fig.2, Fig.3 and Fig.4 of the main article. \nIn SupFig.~\\ref{fig:importations} and \\ref{fig:exportations}, we respectively show the volume of importation and exportation of countries in the world trade network.\n\n\n\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=0.85\\textwidth]{figS1.pdf}\n \\caption{Circles with country flags show country positions on\n the plane of PageRank-CheiRank indexes $(K,K^*)$\n (summation is done over all products) (left panel)\n and on the plane of ImportRank-ExportRank $\\hat{K}$, $\\hat{K}^*$ \n from trade volume (right panel);\n data is shown only for index values less than $61$.\n \n First, second and third row correspond to years 2012, 2014, 2016.}\n \\label{fig:KKstar_2012_2016}\n\\end{figure}\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=0.96\\textwidth]{figS2.pdf}\n \\caption{World map of trade balance of countries \n $B_c = ({P_c}^* - P_c)\/({P_c}^* + P_c)$.\n Left: probabilities are from PageRank and CheiRank vectors;\n Right: probabilities are from the trade volume of Export-Import;\n $B_c$ values are marked by color with the corresponding \n color bar marked by $j$; \n \n \n countries absent in the UN COMTRADE report are marked by black color \n (here and in other SupFigs).\n First, second and third row correspond to year 2012, 2014, 2016.}\n \\label{fig:balance2012_2016}\n\\end{figure}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=0.8\\textwidth]{figS3.pdf}\n \\caption{World map of countries importation volume in dollar.\n Values lower than $10^8$ are clipped for better visibility.\n First, second, third and fourth row correspond to years 2012, 2014, 2016, 2018.\n }\n \\label{fig:importations}\n\\end{figure}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=0.8\\textwidth]{figS4.pdf}\n \\caption{World map of countries exportation volume in dollar.\n Values lower than $10^8$ are clipped for better visibility.\n First, second, third and fourth row correspond to years 2012, 2014, 2016, 2018.\n }\n \\label{fig:exportations}\n\\end{figure}\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=0.96\\textwidth]{figS5.pdf}\n \\caption{Sensitivity of country balance \n $dB_c\/d\\delta_s$ for product $s=3$ (mineral fuels).\n Left: probabilities are from PageRank and CheiRank vectors;\n Right: probabilities are from the trade volume of Export-Import.\n Color bar marked by $j$ \n \n gives sensitivity..\n First, second and third row correspond to year 2012, 2014, 2016.}\n \\label{fig:dbalance_petroleum_2012_2016}\n\\end{figure}\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=0.96\\textwidth]{figS6.pdf}\n \\caption{Same as in Fig.~\\ref{fig:dbalance_petroleum_2012_2016} but\n for product $s=7$ (machinery).}\n \\label{fig:dbalance_machinery_2012_2016}\n\\end{figure}\n\n\n\n\n\n\n\n\n\n\n\n\n\\end{document}","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\nThe spectral heat content represents the total heat \nin a domain $D$ with Dirichlet boundary condition when the initial temperature is 1.\nThe spectral heat content for Brownian motions has been studied extensively.\nThe spectral heat content for isotropic stable processes was first studied in \\cite{Val2017}. \nSince then, considerable progress has been made toward understanding the asymptotic behavior of the spectral heat content for other L\\'evy processes (see \\cite{Val2016, GPS19, P20, PS19}).\n\n\nThe following conjecture about the spectral heat content for isotropic $\\alpha$-stable\nprocesses, $\\alpha\\in (0,2)$, over bounded $C^{1,1}$ open sets (see Section \\ref{section:preliminaries} for the definition of $C^{1,1}$ open sets) was made in \\cite{Val2017}: \nAs $t\\downarrow 0$,\n\\begin{equation}\\label{eqn:Val conj}\nQ_{D}^{(\\alpha)}(t)=\n\\begin{cases}\n|D|-c_{1}|\\partial D|t^{1\/\\alpha} +O(t), \\quad \\alpha\\in (1,2),\\\\\n|D|-c_{2}|\\partial D|t\\ln(1\/t) +O(t), \\quad \\alpha=1,\\\\\n|D|-c_{3}\\text{Per}_{\\alpha}(D)t +o(t), \\quad \\alpha\\in (0,1),\n\\end{cases}\n\\end{equation}\nwhere $c_i, i=1, 2, 3,$ are constants and \n\\begin{equation}\\label{eqn:fractional perimeter}\n\\text{Per}_{\\alpha}(D):=\\int_{D}\\int_{D^{c}}\\frac{A_{d,\\alpha}}{|x-y|^{d+\\alpha}}dydx,\n\\quad \\mbox{ with }A_{d,\\alpha}=\\frac{\\alpha\\Gamma(\\frac{d+\\alpha}{2})}{2^{1-\\alpha}\\pi^{d\/2}\\Gamma(1-\\frac{\\alpha}{2})},\n\\end{equation}\nis the $\\alpha$-fractional perimeter of $D$.\nThis conjecture was resolved in dimension 1 in \\cite{Val2016} (actually, a slightly weaker version, with the error term being $o(t^{1\/\\alpha})$ in the case $\\alpha\\in (1, 2)$ and $o(t\\ln(1\/t))$ in the case $\\alpha=1$, of the conjecture was proved there).\nWe note that in \\cite{Val2017} the author \nconjectured, and also provided strong evidence, that the spectral heat content for isotropic $\\alpha$-stable processes with $\\alpha\\in (0,2)$ must have an asymptotic expansion of the form as \\eqref{eqn:Val conj} for all dimensions $d\\geq2$, \nbut exact expressions for the coefficients $c_{i}$ were not provided.\nThen a two-term asymptotic expansion of the spectral heat content for L\\'evy processes of bounded variation in ${\\mathbb R}^{d}$ was established in \\cite{GPS19}.\nSince $\\alpha$-stable processes are of bounded variation if and only if $\\alpha\\in (0,1)$, \nthe result in \\cite{GPS19} proves \\eqref{eqn:Val conj} for $\\alpha\\in (0,1)$. \nThe purpose of this paper is to resolve the conjecture above for $\\alpha \\in [1,2)$ and $d\\geq 2$. \nIn fact, our result is slightly weaker than \\eqref{eqn:Val conj} since the error term is $o(t^{1\/\\alpha})$ for $\\alpha\\in (1, 2)$ and $o(t\\ln(1\/t))$ for $\\alpha=1$.\nWe also find explicit expressions for the constants $c_{1}$ and $c_{2}$.\nThe main results of this paper are Theorem \\ref{thm:main12} for $\\alpha\\in (1,2)$ and Theorem \\ref{thm:main1} for $\\alpha=1$. \nCombining Theorems \\ref{thm:main12} and \\ref{thm:main1} with \\cite[Corollary 3.5]{GPS19}, the asymptotic behavior of the spectral heat content for isotropic $\\alpha$-stable processes in bounded $C^{1,1}$ open sets $D$ can be stated as follows: \n\n\\begin{thm}\\label{thm:main}\nLet $D$ be a bounded $C^{1,1}$ open set in ${\\mathbb R}^{d}$, $d\\geq 2$ and \nlet\n$$\nf_{\\alpha}(t)=\n\\begin{cases}\nt^{1\/\\alpha} &\\text{ if } \\alpha\\in (1,2),\\\\\nt\\ln(1\/t) &\\text{ if } \\alpha=1,\\\\\nt &\\text{ if } \\alpha\\in (0,1).\n\\end{cases}\n$$\nThen, we have\n\\begin{equation}\\label{eqn:main result}\n\\lim_{t\\to 0}\\frac{|D|-Q_{D}^{(\\alpha)}(t)}{f_{\\alpha}(t)}=\n\\begin{cases}\n{\\mathbb E}[\\overline{Y}_{1}^{(\\alpha)}]|\\partial D| &\\text{if } \\alpha\\in (1,2),\\\\\n\\frac{|\\partial D|}{\\pi} &\\text{if }\\alpha=1,\\\\\n\\rm{Per}_{\\alpha}(D), &\\text{if } \n\\alpha\\in (0,1),\n\\end{cases}\n\\end{equation}\nwhere $\\overline{Y}_{t}^{(\\alpha)}=\\sup_{s\\le t}Y^{(\\alpha)}_s$ stands for the running supremum of a 1 dimensional symmetric $\\alpha$-stable process $Y^{(\\alpha)}_t$ and \n$\\rm{Per}_{\\alpha}(D)$ as defined in \\eqref{eqn:fractional perimeter}.\n\\end{thm}\nWe note that \\eqref{eqn:main result} is exactly the same form as \\cite[Theorem 1.1]{Val2016} if one interprets $|\\partial D|=2$ when $D$ is a bounded open interval in ${\\mathbb R}$, but the proof for $d\\geq 2$ is very different from the one dimensional case and much more challenging.\n\nThe two-term asymptotic expansion of the spectral heat content for Brownian motion was proved in \\cite{BD89}. \nThe crucial ingredient in \\cite{BD89} is the fact that \nindividual components of Brownian motion are independent.\nFor isotropic $\\alpha$-stable processes, $\\alpha\\in (0, 2)$, individual components are not independent and the technique in \\cite{BD89} no longer works.\n When $\\alpha\\in (1,2)$,\nwe establish the lower bound for the heat loss $|D|-Q_{D}^{(\\alpha)}(t)$ by considering the most efficient way of exiting $D$ (see Lemma \\ref{lemma:lb}).\nIn order to establish the upper bound, we approximate the heat loss \n$|D|-Q_{D}^{(\\alpha)}(t)$ by the heat loss of the half-space \n$|D|-Q_{H_{x}}^{(\\alpha)}(t)$ for \n$x$ near $\\partial D$ (see Proposition \\ref{prop:ub12}) \nand show that the approximation error is of order $o(t^{1\/\\alpha})$ in Lemma \\ref{lemma:near boundary}, \nwhich is similar to tools exposed in the trace estimate results in \\cite{BK08, PS14}.\nHowever, these tools do not work when $\\alpha=1$ due to the non-integrability of \n${\\mathbb P}(\\overline{Y}^{(1)}_{1}>u)$ over $(1,\\infty)$, where $ \\overline{Y}^{(1)}_t$ stands for the\nsupremum of the Cauchy process up to time $t$, \nand the proof for $\\alpha=1$ requires new ideas and is considerably more difficult. \n\nIn case of $\\alpha=1$, we prove that the coefficient of the second term \nof the asymptotic expansion of $Q_{D}^{(1)}(t)$ is $-\\frac{|\\partial D|}{\\pi}$, \nwhich is the same as for the regular heat content $H^{(1)}_{D}(t)$ \nwhich represents the total heat in $D$ without the Dirichlet exterior condition \n(see \\eqref{e:rhc} below for the definition of regular heat content and \\cite[Theorem 1.2]{Val2017} \nfor the two-term asymptotic expansion for the regular heat content $H_{D}^{(1)}(t)$). \nIf there is no Dirichlet exterior condition on $D^{c}$, or equivalently if the heat moves freely in and out of $D$,\nthe heat loss of the regular heat content must be smaller than that of the spectral heat content and we obtain the lower bound for free in \\eqref{eqn:Cauchy lb}. The proof for the upper bound is much more demanding. \nThe crucial ingredient for the upper bound is the spectral heat content for subordinate killed Brownian motions in \\cite{PS19}. \nAn isotropic $\\alpha$-stable process $X_{t}^{(\\alpha)}$ can be realized as a subordinate Brownian motion $W_{S_{t}^{(\\alpha\/2)}}$, where $S_{t}^{(\\alpha\/2)}$ is an independent $(\\alpha\/2)$-stable subordinator.\nHence, the spectral heat content $Q_{D}^{(\\alpha)}(t)$ is the spectral heat content for the \\textit{killed subordinate Brownian motion} via the independent $(\\alpha\/2)$-stable subordinator $S_{t}^{(\\alpha\/2)}$. \nWhen one reverses the order of killing and subordination, one obtains the \\textit{subordinate killed Brownian motion}. \nThis is the process obtained by subordinating the killed Brownian motion in $D$ via the independent $(\\alpha\/2)$-stable subordinator $S_{t}^{(\\alpha\/2)}$. \nLet $\\widetilde{Q}_{D}^{(\\alpha)}(t)$ be the spectral heat content of the subordinate killed Brownian motion in $D$\n(see \\eqref{eqn:SKBM} for the precise definition).\nBy construction,\n $\\widetilde{Q}_{D}^{(\\alpha)}(t)$ is always smaller than $Q_{D}^{(\\alpha)}(t)$, \nor equivalently the heat loss $|D|-\\widetilde{Q}_{D}^{(\\alpha)}(t)$ of the subordinate killed Brownian motion provides a natural upper bound for \nthat of the killed subordinate Brownian motion $|D|-Q_{D}^{(\\alpha)}(t)$. \nWe use two independent $\\alpha\/2$ and $\\beta\/2$ stable subordinators \nwith $\\alpha\\beta=2$ \nand consider the $\\alpha$-stable process $X^{(\\alpha)}_{t}=W_{S_{t}^{(\\alpha\/2)}}$ and killed it upon exiting $D$. \nThen we time-change the killed $\\alpha$-stable process by the $(\\beta\/2)$-stable subordinator. \nBy using the heat loss of the resulting process, we obtain in Lemma \n\\ref{lemma:Cauchy ub} that, for $\\alpha\\in (1,2)$, \n$$\n\\limsup_{t\\to 0}\\frac{|D|-Q_{D}^{(1)}(t)}{t\\ln(1\/t)}\\leq\n\\frac{|\\partial D|}{\\pi} +\n\\frac{\\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, Y_{1}^{(\\alpha)}u, Y_{1}^{(\\alpha)}u)-{\\mathbb P}(Y^{(\\alpha)}_{1}\\geq u)\n$$\nand it can be shown that \n${\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u)\\sim {\\mathbb P}(Y^{(\\alpha)}_{1}\\geq u) \\sim cu^{-\\alpha}$ \nas $u\\to \\infty$, \nhence\n$\\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u)du$ and $\\int_{0}^{\\infty}{\\mathbb P}(Y^{(\\alpha)}_{1}\\geq u)du$ \nshould be of order $\\frac{1}{\\alpha-1}$ as $\\alpha\\downarrow 1$.\nWe show that\n$\\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u)du$ and $\\int_{0}^{\\infty}{\\mathbb P}(Y^{(\\alpha)}_{1}\\geq u)du$ \nhave \\textit{the exactly same leading coefficients} and this gives a cancellation of the main terms of order $\\frac{1}{\\alpha-1}$. \nWe still need to show that the sub-leading terms are of order $o(\\alpha-1)$ as $\\alpha\\downarrow 1$ and we show this by establishing uniform heat kernel upper and lower bounds in Lemma \\ref{lemma:uniform ub} \nfor the heat kernel of the supremum process $\\overline{Y}_{t}^{(\\alpha)}$ \nand Lemma \\ref{lemma:uniform lb} for the heat kernel of $Y_{t}^{(\\alpha)}$.\n\n\nThis paper deals with asymptotic behavior of the spectral heat content for isotropic $\\alpha$-stable processes. \nIt is natural and interesting to try to find the asymptotic behavior of the spectral heat content for more general L\\'evy processes. We intend to deal with this topic in a future project.\nIn the recent paper \\cite{P20}, a three-term asymptotic expansion of the spectral heat content of 1 dimensional symmetric $\\alpha$-stable processes, $\\alpha\\in [1,2)$, was established.\nWe believe that a similar result should hold true for $d\\geq 2$.\n\nThe organization of this paper is as follows. \nIn Section \\ref{section:preliminaries}, we introduce the setup. In Section \\ref{section:alpha12},\nwe deal with the case $\\alpha\\in (1,2)$ and the main result of that section is Theorem \\ref{thm:main12}. \nThe case $\\alpha=1$ is dealt with in Section \\ref{section:alpha1} and the main result there is Theorem \\ref{thm:main1}. \nIn this paper, we use $c_i$ to denote constants whose values are unimportant and may change from one appearance to another.\nThe notation ${\\mathbb P}_{x}$ stands for the law of the underlying processes started at $x\\in {\\mathbb R}$, and ${\\mathbb E}_{x}$ stands for expectation with respect to ${\\mathbb P}_{x}$. For simplicity, we use ${\\mathbb P}={\\mathbb P}_{0}$ and ${\\mathbb E}={\\mathbb E}_{0}$.\n\n\n\\section{Preliminaries}\\label{section:preliminaries}\n\nIn this paper, unless explicitly stated otherwise, we assume $d\\geq 2$.\nLet $X_{t}^{(\\alpha)}$, $\\alpha\\in (0, 2]$, be an isotropic $\\alpha$-stable process with \n$$\n{\\mathbb E}[e^{i\\xi X^{(\\alpha)}_{t}}]=e^{-t|\\xi|^{\\alpha}}, \\quad \\xi \\in {\\mathbb R}^{d}, \\alpha\\in (0,2].\n$$\n$X_{t}^{(2)}$ is a Brownian motion $W_t$ with transition density given by $(4\\pi t)^{-d\/2}e^{-\\frac{|x|^2}{4t}}$. \nLet $S_{t}^{(\\alpha\/2)}$, $\\alpha\\in (0,2)$, be an $(\\alpha\/2)$-stable subordinator with\n$$\n{\\mathbb E}[e^{-\\lambda S_{t}^{(\\alpha\/2)}}]=e^{-t \\lambda^{\\alpha\/2}}, \\quad \\lambda>0.\n$$\nAssume that $S_{t}^{(\\alpha\/2)}$ is independent of the Brownian motion $W_t$. \nThen, the subordinate Brownian motion $W_{S_{t}^{(\\alpha\/2)}}$ is a realization of the process $X_{t}^{(\\alpha)}$.\nWe will reserve $Y_{t}^{(\\alpha)}$ for the 1 dimensional symmetric $\\alpha$-stable process.\nWe define the running supremum process $\\overline{Y}^{(\\alpha)}_{t}$ of $Y_{t}^{(\\alpha)}$ by\n\\begin{equation}\\label{eqn:running}\n\\overline{Y}^{(\\alpha)}_{t}=\\sup\\{Y^{(\\alpha)}_{u} : 0\\leq u\\leq t\\}.\n\\end{equation}\n\nRecall that an open set $D$ in ${\\mathbb R}^d$ is said to be a $C^{1, 1}$ open set if \nthere exist a localization radius $R_0>0$ and a constant $\\Lambda_0>0$ such that, \nfor every $z\\in \\partial D$, there exist a $C^{1, 1}$ function $\\phi=\\phi_z: {\\mathbb R}^{d-1}\\to {\\mathbb R}$ satisfying $\\phi(0)=0$, $\\nabla \\phi(0)=(0, \\cdots, 0)$, $\\|\\nabla\\phi\\|_\\infty\\le \\Lambda_0$,\n$|\\nabla \\phi(x)-\\nabla \\phi(y)|\\le \\Lambda_0 |x-y|$ and an orthonormal coordinate system $CS_z: y=(\\widetilde y, y_d)$\nwith origin at $z$ such that\n$$\nB(z, R_0)\\cap D=B(z, R_0)\\cap \\{ y=(\\widetilde y, y_d) \\mbox{ in } CS_z: y_d>\\phi(\\widetilde y)\\}.\n$$\nThe pair $(R_0, \\Lambda_0)$ is called the $C^{1,1}$ characteristics of the $C^{1, 1}$ open set $D$. \nIt is well known that any $C^{1, 1}$ open set $D$ in ${\\mathbb R}^d$ satisfies the uniform interior and exterior $R$-ball condition: for any $z\\in\\partial D$, there exist balls $B_1$ and $B_2$ of radii $R$ with $B_1\\subset D$, $B_2\\subset{\\mathbb R}^d\\setminus\\overline{D}$ and $\\partial B_1\\cap\\partial B_2=\\{z\\}$.\n\nWe recall from \\cite{BD89} a useful fact about open sets $D$ satisfying the uniform interior and exterior $R$-ball condition.\nLet $D_{q}=\\{x\\in D : \\text{dist}(x, \\partial D) >q\\}$.\nWe will use $\\partial D_{q}$ denote the portion of the boundary of $D_q$ contained in $D$, that is, $\\partial D_{q}=\\{x\\in D: \\text{dist}(x, \\partial D)=q\\}$.\nIt follows from \\cite[Lemma 6.7]{BD89}\nthat\n\\begin{equation}\\label{eqn:stability}\n|\\partial D|\\left(\\frac{R-q}{R}\\right)^{d-1}\\leq |\\partial D_{q}|\\leq |\\partial D|\\left(\\frac{R}{R-q}\\right)^{d-1}, \\quad 0< q0: X^{(\\alpha)}_{t}\\notin D\\}$ be the first exit time of $X^{(\\alpha)}_{t}$ from $D$. \nThe spectral heat content of $D$ for $X^{(\\alpha)}_{t}$ is defined to be\n$$\nQ_{D}^{(\\alpha)}(t)=\\int_{D}{\\mathbb P}_{x}(\\tau_{D}^{(\\alpha)}>t)dx.\n$$\nThe (regular) heat content $H_{D}^{(\\alpha)}(t)$ of $D$ for $X^{(\\alpha)}_{t}$ is defined to be\n\\begin{equation}\\label{e:rhc}\nH_{D}^{(\\alpha)}(t)=\\int_{D}{\\mathbb P}_{x}(X_{t}^{(\\alpha)}\\in D)dx.\n\\end{equation}\nThe spectral heat content $\\widetilde{Q}_{D}^{(\\alpha)}(t)$ of $D$ for the subordinate killed Brownian motion is defined as \n\\begin{equation}\\label{eqn:SKBM}\n\\widetilde{Q}_{D}^{(\\alpha)}(t)=\\int_{D}{\\mathbb P}_{x}(\\tau_{D}^{(2)}>S_{t}^{(\\alpha\/2)})dx.\n\\end{equation}\nNote that we always have \n\\begin{equation}\\label{eqn:two heat contents}\n\\{\\tau_{D}^{(2)}>S_{t}^{(\\alpha\/2)}\\} \\subset \\{\\tau_{D}^{(\\alpha)}>t\\}\\subset \\{X_{t}^{(\\alpha)}\\in D\\}.\n\\end{equation}\nWe sometimes use the terminology \\textit{heat loss} and this will mean either\n$$\n|D|-Q_{D}^{(\\alpha)}(t)=\\int_{D}{\\mathbb P}_{x}(\\tau_{D}^{(\\alpha)}\\leq t)dx,\n$$\nor \n$$\n|D|-\\widetilde{Q}_{D}^{(\\alpha)}(t)=\\int_{D}{\\mathbb P}_{x}(\\tau_{D}^{(2)}\\leq S_{t}^{(\\alpha\/2)})dx,\n$$\ndepending on which process we are dealing with. \nIntuitively, these quantities represent the total heat loss caused by heat particles jumping out of $D$ up to time $t$.\nFrom \\eqref{eqn:two heat contents}, we have \n$$\n|D|- H_{D}^{(\\alpha)}(t)\\leq |D|-Q_{D}^{(\\alpha)}(t) \\leq |D|-\\widetilde{Q}_{D}^{(\\alpha)}(t), \\quad t>0.\n$$\n\n\n\n\\section{The case $\\alpha\\in (1,2)$}\\label{section:alpha12}\nThroughout this section, we assume $\\alpha\\in(1,2)$.\nFor any $x\\in D$, we use $\\delta_D(x)$ to denote the distance between $x$ and $\\partial D$.\nWe start with a lower bound. \nRecall that $D_{q}=\\{x\\in D : \\text{dist}(x, \\partial D) >q\\}$.\n\n\n\n\\begin{lemma}\\label{lemma:lb}\nLet $D$ be a bounded $C^{1,1}$ open set in ${\\mathbb R}^{d}$. Then, we have\n$$\n\\liminf_{t\\to 0}\\frac{|D|-Q^{(\\alpha)}_{D}(t)}{t^{1\/\\alpha}}\\geq |\\partial D|\n{\\mathbb E}[\\overline{Y}^{(\\alpha)}_{1}],\n$$\nwhere $\\overline{Y}_{t}^{(\\alpha)}$ is defined in \\eqref{eqn:running}.\n\\end{lemma}\n\\noindent{\\bf Proof.} \nLet $D$ satisfy the uniform interior and exterior $R$-ball condition.\nFix $a\\leq R\/2$. \nFor $x\\in D\\setminus D_{a}$, let \n$z_x\\in \\partial D$ be such that $|x-z_x|=\\delta_D(x)$ and let ${\\bf n}_{z_x}$ be the outward unit normal vector to the boundary $\\partial D$ at the point $z_x$.\n\nFor $X$ starting from $x$, we define \n$Y^{(\\alpha)}_{t}:=(X_{t}^{(\\alpha)}-x)\\cdot {\\bf n}_{z_x}$, where $\\cdot$ stands for usual scalar product in ${\\mathbb R}^{d}$. Obviously, $Y^{(\\alpha)}_{0}=0$.\nNote that the characteristic exponent of $Y^{(\\alpha)}_{t}$ is given by \n$$\n{\\mathbb E}[e^{i\\eta Y^{(\\alpha)}_{t}}]=\n{\\mathbb E}_x[e^{i\\eta (X_{t}^{(\\alpha)}-x)\\cdot{\\bf n}_{z_x}}]\n=e^{-t|\\eta|^{\\alpha}}, \n\\quad \\eta \\in {\\mathbb R},\n$$\nand this shows that $Y^{(\\alpha)}_{t}$ is a one dimensional stable process starting from 0.\nLet $r\\le a$ and $x\\in \\partial D_{r}$, and let $H_{x}$ be the half-space containing the interior $R$-ball at the point $z_x$ and tangent to $\\partial D$ at $z_x$.\nWhen the process $X$ starts from $x$, we have\n$$\n\\{\\overline{Y}_{t}^{(\\alpha)} >r\\}=\\{\\tau_{H_x}^{(\\alpha)}\\leq t\\}\\subset \\{\\tau_{D}^{(\\alpha)} \\leq t\\} \n$$\nand this implies \n$$\n{\\mathbb P}(\\overline{Y}_{t}^{(\\alpha)} >r)\\leq {\\mathbb P}_{x}(\\tau_{D}^{(\\alpha)} \\leq t).\n$$\nHence, by the coarea formula\nand \\eqref{eqn:stability}, we have\n\\begin{align*}\n&|D|-Q^{(\\alpha)}_{D}(t)=\\int_{D}{\\mathbb P}_{x}(\\tau^{(\\alpha)}_{D}\\leq t)dx\n\\geq \\int_{D\\setminus D_{a}}{\\mathbb P}_{x}(\\tau^{(\\alpha)}_{D}\\leq t)dx\\\\\n&=\\int_{0}^{a}\\int_{\\partial D_{r}}{\\mathbb P}_{x}(\\tau^{(\\alpha)}_{D}\\leq t)dSdr\n\\geq \\int_{0}^{a}\\int_{\\partial D_{r}}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{t}>r)dSdr\\\\\n&=\\int_{0}^{a}|\\partial D_{r}|{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{t}>r)dr\n\\geq |\\partial D|\\left(\\frac{R-a}{R}\\right)^{d-1}\\int_{0}^{a}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{t}>r)dr,\n\\end{align*}\nwhere $dS$ represents the surface measure on $\\partial D_{r}$.\nNow it follows from the scaling property of $Y^{(\\alpha)}$ and the change of variables $t^{-1\/\\alpha}r=s$ that\n\\begin{align*}\n&\\int_{0}^{a}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{t}>r)dr\n=\\int_{0}^{a}{\\mathbb P}(t^{1\/\\alpha}\\overline{Y}^{(\\alpha)}_{1}>r)dr\n=t^{1\/\\alpha}\\int_{0}^{t^{-1\/\\alpha}a}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>s)ds.\n\\end{align*}\nHence, we have\n$$\n\\frac{|D|-Q^{(\\alpha)}_{D}(t)}{t^{1\/\\alpha}}\\geq |\\partial D|\\left(\\frac{R-a}{R}\\right)^{d-1}\\int_{0}^{t^{-1\/\\alpha}a}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>s)ds\n$$\nand by taking $\\liminf$ we obtain\n$$\n\\liminf_{t\\to 0}\\frac{|D|-Q^{(\\alpha)}_{D}(t)}{t^{1\/\\alpha}}\\geq |\\partial D|\\left(\\frac{R-a}{R}\\right)^{d-1}{\\mathbb E}[\\overline{Y}^{(\\alpha)}_1].\n$$\nSince $a>0$ is arbitrary and \n${\\mathbb E}[\\overline{Y}^{(\\alpha)}_{1}]<\\infty$ \nwhen $\\alpha\\in (1,2)$, the conclusion of the theorem is true.\n{\\hfill $\\Box$ \\bigskip}\n\n\n\n\nNow we proceed to prove the opposite inequality of Lemma \\ref{lemma:lb}. \nFor a bounded $C^{1,1}$ open set $D$ satisfying the uniform \ninterior and exterior $R$-ball condition\nand $x\\in D$ with $\\delta_{D}(x)t\\}\\nonumber\\\\\n&\\subset&\\{\\tau_{H_{x}}^{(\\alpha)}\\leq t\\} \\cup \\{\\tau_{D}^{(\\alpha)}\\leq t <\\tau_{H_{x}}^{(\\alpha)}\\}.\n\\end{eqnarray}\nHence, we have\n\\begin{eqnarray}\\label{eqn:ub cases}\n&&|D|-Q^{(\\alpha)}_{D}(t)=\\int_{D}{\\mathbb P}_{x}(\\tau^{(\\alpha)}_{D}\\leq t)dx=\\int_{D_{R\/2}}{\\mathbb P}_{x}(\\tau_{D}^{(\\alpha)}\\leq t)dx +\\int_{D\\setminus D_{R\/2}}{\\mathbb P}_{x}(\\tau_{D}^{(\\alpha)}\\leq t)dx\\nonumber\\\\\n&\\leq &\\int_{D_{R\/2}}{\\mathbb P}_{x}(\\tau^{(\\alpha)}_{D}\\leq t)dx+\\int_{D\\setminus D_{R\/2}}{\\mathbb P}_{x}(\\tau^{(\\alpha)}_{H_{x}}\\leq t)dx +\\int_{D\\setminus D_{R\/2}}{\\mathbb P}_{x}(\\tau^{(\\alpha)}_{D}\\leq t<\\tau^{(\\alpha)}_{H_{x}})dx.\n\\end{eqnarray}\n\nWe deal with the first expression of \\eqref{eqn:ub cases} first.\nRecall the following facts from \\cite[(2.2)]{Val2017} and \\cite[Corollary 6.4]{BD89}:\n\\begin{equation}\\label{eqn:exp moment}\n{\\mathbb E}\\left[\\exp\\left(-\\frac{\\kappa^{2}}{S_{1}^{(\\alpha\/2)}}\\right)\\right]\\leq c_{\\alpha}\\kappa^{-\\alpha}\n\\end{equation}\nand\n\\begin{equation}\\label{eqn:BM exit}\n{\\mathbb P}\\left(\\tau^{(2)}_{B(0,R)}0\\}$.\n\\end{prop}\n\\noindent{\\bf Proof.} \nFor $\\delta_{D}(x)r\\}$, \nand this implies\n${\\mathbb P}_{(\\widetilde 0, r)}(\\tau^{(\\alpha)}_{H}\\leq 1)=\n{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>r)$. \nHence, we have\n$$\n\\int_{0}^{\\infty}{\\mathbb P}_{(\\widetilde 0, v)}(\\tau^{(\\alpha)}\\leq 1)dv\n=\\int_{0}^{\\infty}{\\mathbb P}(\n\\overline{Y}^{(\\alpha)}_{1}\\geq v)dv={\\mathbb E}[\\overline{Y}^{(\\alpha)}_{1}].\n$$\nBy \\cite[Corollary 2.1 (ii)]{Val2016}, ${\\mathbb E}[\\overline{Y}^{(\\alpha)}_{1}]<\\infty$ if $\\alpha\\in (1,2)$.\nHence, it follows from the Lebesgue dominated convergence theorem that\n$$\n\\lim_{t\\to 0}\\frac{\\int_{D\\setminus D_{R\/2}} {\\mathbb P}_{x}(\\tau^{(\\alpha)}_{H_{x}} \\leq t)dx}{t^{1\/\\alpha}}=|\\partial D|\\int_{0}^{\\infty}\n{\\mathbb P}_{(\\widetilde 0, r)}(\\tau^{(\\alpha)}_{H}\\leq 1)dr=\n|\\partial D|{\\mathbb E}[\\overline{Y}^{(\\alpha)}_{1}].\n$$\n{\\hfill $\\Box$ \\bigskip}\n\n\n\nFinally, we estimate the last expression of \\eqref{eqn:ub cases}.\n\\begin{lemma}\\label{lemma:near boundary}\nLet $D\\subset {\\mathbb R}^{d}$, $d\\geq 2$, be a bounded $C^{1,1}$ \nopen set. \nThen, we have\n$$\n\\lim_{t\\to 0}\\frac{\\int_{D\\setminus D_{R\/2}}{\\mathbb P}_{x}\\left(\\tau_{D}^{(\\alpha)}\\leq t< \\tau^{(\\alpha)}_{H_{x}}\\right)dx}{t^{1\/\\alpha}}=0.\n$$\n\\end{lemma}\n\\noindent{\\bf Proof.} \nBy rotational invariance, we have\n\\begin{align}\\label{eqn:approx1}\n&\\int_{D\\setminus D_{R\/2}}{\\mathbb P}_{x}\\left(\\tau_{D}^{(\\alpha)}\\leq t< \\tau^{(\\alpha)}_{H_{x}}\\right)dx\n\\leq\\int_{D\\setminus D_{R\/2}}{\\mathbb P}_{x}\\left(\\tau^{(\\alpha)}_{B_{x}}\\leq t< \\tau^{(\\alpha)}_{H_{x}}\\right)dx\\nonumber\\\\\n&\\leq \\int_{0}^{R\/2}|\\partial D_{u}|{\\mathbb P}_{(\\widetilde 0, u)}\n\\left(\\tau^{(\\alpha)}_{B((\\widetilde 0, R), R)}\n\\leq t< \\tau^{(\\alpha)}_{H}\\right)du,\n\\end{align}\nwhere $H=\\{x=(x_{1},\\cdots ,x_{d}) : x_{d}>0\\}$.\nIt follows from \\cite[Lemma 6.7]{BD89} that for $uu, Y_{1}^{(\\alpha)}u, Y_{1}^{(\\alpha)}b, Y_{t}^{(\\alpha)}\\in (a,b))dx +\\int_{a}^{b}{\\mathbb P}_{x}(\\underline{Y}^{(\\alpha)}_{t}b, \\underline{Y}^{(\\alpha)}_{t}b, \\underline{Y}^{(\\alpha)}_{t}b, Y_{t}^{(\\alpha)}\\in (a,b))dx\\\\\n&=&t^{1\/\\alpha}\\int_{0}^{(b-a)t^{-1\/\\alpha}}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, -(b-a)t^{-1\/\\alpha}+u u, -(b-a)t^{-1\/\\alpha}+u u)\n$$\nand\n$$\n\\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u)du={\\mathbb E}[\\overline{Y}^{(\\alpha)}_{1}]<\\infty.\n$$\nHence, it follows from the Lebesgue dominated convergence theorem that the limit is\n\\begin{align*}\n&\\lim_{t\\to 0}\\frac{\\int_{a}^{b}{\\mathbb P}_{x}(\\tau_{(a,b)}^{(\\alpha)}\\leq t, Y_{t}^{(\\alpha)}\\in (a,b))dx}{t^{1\/\\alpha}}\\\\\n=&\\lim_{t\\to 0}\\frac{t^{1\/\\alpha}\\int_{0}^{(b-a)t^{-1\/\\alpha}}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, -(b-a)t^{-1\/\\alpha}+u u, Y_{1}^{(\\alpha)}0$ and $\\alpha\\in (0,2)$, we have\n$$\n\\lim_{t\\to 0}\\frac{{\\mathbb E}[\\left(S_{1}^{(\\alpha\/2)}\\right)^{\\alpha\/2}, 0u, Y_{1}^{(\\alpha)}T_{t}^{(\\beta\/2)})dx.\n$$\n\nSince $\\{\\tau_{D}^{(1)} \\leq t\\}\\subset \\{\\tau_{D}^{(\\alpha)} \\leq T_{t}^{(\\beta\/2)}\\}$, we have $|D|-Q_{D}^{(1)}(t)\\leq |D|-\\widetilde{Q}^{(\\alpha,\\beta)}_{D}(t)$. \nWe will show that \n\\begin{equation}\\label{eqn:Cauchy aux7}\n\\limsup_{t\\to 0}\\frac{|D|-\\widetilde{Q}^{(\\alpha,\\beta)}_{D}(t)}{t\\ln(1\/t)}\\leq \n\\frac{|\\partial D|}{\\pi} +\n\\frac{\\int_{0}^{\\infty}\n{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, Y_{1}^{(\\alpha)}0$,\nthere exists $\\delta=\\delta(\\varepsilon)$ such that for all $u\\leq \\delta$,\n\\begin{eqnarray*}\n\\frac{|D|-Q_{D}^{(\\alpha)}(u)}{u^{1\/\\alpha}} \n\\leq \\left(\\frac{\\Gamma(1-\\frac{1}{\\alpha})}{\\pi} + \\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, Y_{1}^{(\\alpha)}u, Y_{1}^{(\\alpha)}0$ is arbitrary, this establishes \\eqref{eqn:Cauchy aux7}.\n{\\hfill $\\Box$ \\bigskip}\n\n\nWe now recall the definition for the double gamma function from \\cite[(4.4)]{Ku11}. \nFor $z\\in \\mathbb{C}$ and $\\tau \\in \\mathbb{C}$ with $|\\text{arg}(\\tau)|<\\pi$, we define \n$$\nG(z;\\tau)=\\frac{z}{\\tau}e^{\\frac{az}{\\tau} + \\frac{bz^2}{2\\tau}}\n{\\textstyle \\prod}_{m\\geq 0}{\\textstyle \\prod}^{'}_{n\\geq 0}\\left( 1+\\frac{z}{m\\tau +n}\\right)\\exp\\left(-\\frac{z}{m\\tau +n} +\\frac{z^2}{2(m\\tau +n)^2}\\right),\n$$\nwhere the prime in the second product means that the term corresponding to $m=n=0$ is omitted. \n\n\\begin{lemma}\\label{lemma:double gamma bounded}\nLet $K$ be a compact set in $\\mathbb{C}$. \nThen there exists a constant $c=c(K)>0$ such that \n$$\n\\left| G(z;\\tau) \\right| \\leq c \\text{ for all } z\\in K \\text{ and } \\tau\\in [1,2].\n$$\n\\end{lemma}\n\\noindent{\\bf Proof.} \nOn the compact set $K\\times [1, 2]$, $\\frac{z}{\\tau}e^{\\frac{az}{\\tau} + \\frac{bz^2}{2\\tau}}$ is continuous and hence bounded. \nWe only need to prove that the double infinite product in $G(z, \\tau)$ is bounded. \nRecall the canonical factor $E_{2}(z)$ (see \\cite[p. 145]{SS complex}):\n$$\nE_{2}(z)=(1-z)e^{z+\\frac{z^2}{2}}.\n$$\nThe double infinite product in $G(z;\\tau)$ can be written as\n$$\n {\\textstyle \\prod}_{m\\geq 0}{\\textstyle \\prod}^{'}_{n\\geq 0}\\left( 1+\\frac{z}{m\\tau +n}\\right)\\exp\\left(-\\frac{z}{m\\tau +n} +\\frac{z^2}{2(m\\tau +n)^2}\\right)\n = {\\textstyle \\prod}_{m\\geq 0}{\\textstyle \\prod}^{'}_{n\\geq 0}E_{2}(-\\frac{z}{m\\tau +n}).\n$$\nIt follows from \\cite[Lemma 4.2]{SS complex} for $|w|\\leq 1\/2$, $|1-E_{2}(w)|\\leq c_{1}|w|^{3}$ for some constant $c_{1}$.\nSince $|-\\frac{z}{m\\tau +n}| \\leq \\frac{|z|}{m+n}$, \nall but finitely many terms of the form $|-\\frac{z}{m\\tau +n}|$ in the double infinite product are less than $\\frac12$, and since each of these finitely many terms are bounded on $K\\times [1, 2]$, we may assume $|-\\frac{z}{m\\tau +n}|\\leq \\frac12$ on $K\\times [1, 2]$ for all $n,m$.\n\nNow it follows from \\cite[Lemma 4.2]{SS complex} for all \n$(z, \\tau)\\in K\\times [1, 2]$ that\n$$\n\\left| 1-E_{2}(-\\frac{z}{m\\tau +n})\\right| \\leq c_{1}\\left |\\frac{z}{m\\tau +n}\\right|^{3}\\leq \\frac{c_{2}}{(m+n)^{3}} \\text{ for all } n,m.\n$$\nIt is easy to see that \n$$\n\\sum_{m\\geq 0, n\\geq 0, (m,n)\\neq (0, 0)}\\frac{1}{(m+n)^{3}}<\\infty.\n$$\nNote that for any $|w|\\leq \\frac{1}{2}$ we have $\\ln |1-w|\\leq 2|w|$.\nThere exist positive integers $N_1$ and $M_1$ such that $|1-E_{2}(-\\frac{z}{m\\tau +n})|\\leq \\frac{1}{2}$ on $K\\times [1, 2]$ for all $n\\geq N_{1}$ and $m\\geq M_{1}$.\nHence, for any $M, N \\in {\\mathbb N}$ large we have \n\\begin{align*}\n&\\left|{\\textstyle \\prod}_{M_{1}\\leq m \\leq M}{\\textstyle \\prod}_{N_{1}\\leq n\\leq N} E_{2}(-\\frac{z}{m\\tau +n}) \\right|\n={\\textstyle \\prod}_{M_{1}\\leq m \\leq M}{\\textstyle \\prod}_{N_{1}\\leq n\\leq N}\\left|1- (1- E_{2}(-\\frac{z}{m\\tau +n}) )\\right|\\\\\n&={\\textstyle \\prod}_{M_{1}\\leq m \\leq M}{\\textstyle \\prod}_{N_{1}\\leq n\\leq N}\\exp \\left( \\ln \\left| (1- (1- E_{2}(-\\frac{z}{m\\tau +n}) )) \\right| \\right)\\\\\n&= \\exp \\left( \\sum_{M_{1}\\leq m \\leq M,N_{1}\\leq n\\leq N} \\ln \\left|(1- (1- E_{2}(-\\frac{z}{m\\tau +n}) )) \\right| \\right)\\\\\n&\\leq \\exp \\left( 2\\sum_{M_{1}\\leq m \\leq M,N_{1}\\leq n\\leq N} \\left| 1- E_{2}(-\\frac{z}{m\\tau +n}) \\right| \\right)\\leq \\exp \\left( 2c_2\\sum_{M_{1}\\leq m \\leq M,N_{1}\\leq n\\leq N} \\frac{1}{(m+n)^{3}} \\right)\\\\\n&\\leq \\exp \\left( 2c_2\\sum_{M_{1}\\leq m, N_{1}\\leq n} \\frac{1}{(m+n)^{3}} \\right).\n\\end{align*}\nBy letting $M, N\\to \\infty$ we see that the double infinite product is bounded on \n$K\\times [1, 2]$.\n{\\hfill $\\Box$ \\bigskip}\n\nIt follows from \\cite[Proposition VIII.1-4]{Ber} that there exists a constant $C>0$ such that\n\\begin{equation}\\label{eqn:same asymptotic}\n{\\mathbb P}(\\overline{Y}^{(\\alpha)}_1>u)\\sim{\\mathbb P}(Y^{(\\alpha)}_1>u)\\sim Cu^{-\\alpha} \n\\text{ as } u\\to \\infty.\n\\end{equation}\nLet $\\rho:={\\mathbb P}(Y_1^{(\\alpha)}>0)$. We say $Y^{(\\alpha)}\\in C_{k,l}$ if\n\\begin{align}\\label{eqn:ckl}\n&(\\alpha, \\rho)\\in \\{\\alpha\\in (0,1), \\rho\\in(0,1)\\} \\cup \\{\\alpha=1, \\rho=1\/2\\}\\cup \\{\\alpha\\in (1,2), \\rho\\in [1-\\frac{1}{\\alpha}, \\frac{1}{\\alpha}]\\} \\text{ and}\\nonumber\\\\\n&\\frac{1}{2}+k=\\frac{l}{\\alpha}, \\text{ or equivalently } \\alpha=\\frac{2l}{1+2k} \\text{ for some } k,l\\in {\\mathbb N}.\n\\end{align}\nNote that this condition already appeared in \\cite[\\text{Definition 1}]{Ku11}. \n\n\\begin{lemma}\\label{lemma:uniform ub}\nSuppose that $Y^{(\\alpha)}\\in C_{k,k+1}$\nfor some $k\\in \\mathbb{N}$.\nThen $\\overline{Y}_1^{(\\alpha)}$ has a density $\\overline{p}^{(\\alpha)}(x)$ and there exists \na constant $A>0$, \nindependent of all $\\alpha\\in(1,2)$ satisfying $Y_{t}^{(\\alpha)}\\in C_{k,k+1}$, such that\n$$\n\\overline{p}^{(\\alpha)}(x)\\leq C\\alpha x^{-1-\\alpha} +Ax^{-3}\n$$\nfor all $x>0$, where $C$ is the constant in \\eqref{eqn:same asymptotic}.\n\\end{lemma}\n\\noindent{\\bf Proof.} \nThe proof is similar to that of \\cite[Theorem 9]{Ku11} with a focus on establishing \na uniform constant $A$. \nIt follows from the proof of \\cite[Theorem 9]{Ku11} that $\\overline{Y}_1^{(\\alpha)}$ has a density $\\overline{p}^{(\\alpha)}(x)$ given by the inverse Mellin transform\n$$\n\\overline{p}^{(\\alpha)}(x)=\\frac{1}{2\\pi i}\\int_{1+i{\\mathbb R}}M(s,\\alpha)x^{-s}ds,\n$$\nwhere $M(s,\\alpha)$ is the Mellin transform given by\n$$\nM(s,\\alpha)={\\mathbb E}[(\\overline{Y}^{(\\alpha)}_{1})^{s-1}], \\quad s\\in \\mathbb{C}.\n$$\nWe remark here that we write the Mellin transform as $M(s,\\alpha)$ instead of $M(s)$ to emphasize its dependence on $\\alpha$.\nIt follows from \\cite[Lemma 2]{Ku11} that $M(s,\\alpha)$ can be extended to a meromorphic function on $\\mathbb{C}$ whose simples poles are at $s_{m,n}:=m+\\alpha n$, where $m\\leq 1-(k+1)$ and $n\\in \\{0,1,\\cdots, k\\}$, or $m\\geq 1$ and $n\\in \\{1,2,\\cdots, k\\}$ with residues\n$$\n\\text{Res}(M(s,\\alpha), s_{m,n})=c_{m-1,n}^{+},\n$$\nwhere $c^{+}_{m-1,n}$ is the constant defined in \\cite[(7.5)]{Ku11}.\n\nSince $Y^{(\\alpha)}\\in C_{k,k+1}$, \nit follows from \\eqref{eqn:ckl} that $\\alpha=\\frac{2+2k}{1+2k}\\in (1,2)$.\nIf $m\\geq 1$ and $n\\in \\{1,\\cdots, k\\}$ then $s_{m,n}=m+\\alpha n$ has a nonempty intersection with $[1,3]$ if and only if $m=n=1$. In particular, $s_{1,1}=1+\\alpha$. \nIf $m\\leq 1-(k+1)=-k$ and $n\\in \\{0,\\cdots, k\\}$, we have \n$$\ns_{m,n}=m+\\alpha n \\leq -k +\\frac{2+2k}{1+2k}k=\\frac{k}{1+2k} <\\frac{1}{2}.\n$$\nHence, the only simple pole of $M(s,\\alpha)$ with $\\text{Re}(s)\\in [1,3]$ is $s_{1,1}=1+\\alpha$.\n\n\nWe note here that although it is written as $y\\to \\infty$ in \\cite[Lemma 3]{Ku11}, it is actually true for all $|y|\\to \\infty$ by checking the proof there as \\cite[(7.15)]{Ku11} is true for all $z\\to\\infty$ with $|\\text{arg}(z)|<\\pi$.\nHence, by taking a rectangle with vertices at $1\\pm Pi$ and $3\\pm Pi$ with the help of the residue theorem, then letting $P\\to \\infty$ and using \\cite[Lemma 3]{Ku11}, we get\n\\begin{eqnarray*}\n&&\\overline{p}^{(\\alpha)}(x)=\\frac{1}{2\\pi i}\\int_{1+i{\\mathbb R}}M(s,\\alpha)x^{-s}ds \\\\\n&=&\\text{Res}(M(s,\\alpha), s_{1,1}^{+})x^{-1-\\alpha} +\\frac{1}{2\\pi i}\\int_{3+i{\\mathbb R}}M(s,\\alpha)x^{-s}ds\\\\\n&=&c_{0,1}^{+}x^{-1-\\alpha}+\\frac{1}{2\\pi i}\\int_{3+i{\\mathbb R}}M(s,\\alpha)x^{-s}ds,\n\\end{eqnarray*}\nwhere $\\text{Res}(M(s,\\alpha), s_{1,1}^{+})$ is the residue of $M(s,\\alpha)$ at $s_{1,1}^{+}$.\n\nWe claim that the constant $c_{0,1}^{+}$ must be $C\\alpha$, where $C$ is from \\eqref{eqn:same asymptotic} as we will show that the reminder is $O(x^{-3})$, which will in turn imply that there exist constants $c_{1}, c_{2}\\in {\\mathbb R}$ such that \n$$\nc_{0,1}^{+}x^{-1-\\alpha} +c_1x^{-3} \\leq \\overline{p}^{(\\alpha)}(x) \\leq c_{0,1}^{+}x^{-1-\\alpha} +c_2x^{-3}\n$$\nfor all sufficiently large $x$.\nBy integrating on $(u,\\infty)$ we obtain \n$$\n\\frac{c_{0,1}^{+}}{\\alpha}u^{-\\alpha} +\\frac{c_1}{2} u^{-2}\\leq \n{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u) \n\\leq \\frac{c_{0,1}^{+}}{\\alpha}u^{-\\alpha} +\\frac{c_2}{2}u^{-2},\n$$\nfor all sufficiently large $u>0$.\nComparing the equation above with \\eqref{eqn:same asymptotic} we conclude that \n$\\alpha^{-1}c_{0,1}^{+}=C$.\nThis shows that the leading term of $\\overline{p}^{(\\alpha)}(x)$ is $C\\alpha x^{-1-\\alpha}$.\n\nHence, the proof will be complete once we show that \n\\begin{equation}\\label{eqn:Mellin error}\n\\left| \\frac{1}{2\\pi i}\\int_{3+i{\\mathbb R}}M(s,\\alpha)x^{-s}ds\\right| \n\\leq Ax^{-3},\n\\end{equation}\nwhere the constant $A$ is independent of all $\\alpha\\in (1,2)$ satisfying \n$Y^{(\\alpha)}\\in C_{k,k+1}$.\n\nFrom \\cite[Lemma 3]{Ku11} we know that \n\\begin{equation}\\label{eqn:Mellin}\n\\ln|M(3+iy,\\alpha)|=-\\frac{\\pi |y|}{2\\alpha} +o(y) \\text{ as } |y|\\to\\infty.\n\\end{equation}\nThe proof of \\eqref{eqn:Mellin} in \\cite[Lemma 3]{Ku11} depends on \\cite[(4.5)]{BK97}, which is a uniform estimate of the double gamma functions. Hence, the error term in \\eqref{eqn:Mellin} is uniform for all $\\alpha\\in (1,2)$ and \nthere exists $N>0$, independent of $\\alpha\\in (1,2)$, such that \n$$\n\\ln|M(3+iy,\\alpha)|\\leq -\\frac{\\pi |y|}{5} \\text{ for all } |y|\\geq N \\text{ and } \\alpha\\in(1,2).\n$$\nHence, we have \n\\begin{equation}\\label{eqn:Mellin ub}\n|M(3+iy, \\alpha)|\\leq e^{-\\frac{\\pi|y|}{5}} \\text{ for all } |y|\\geq N \\text{ and } \\alpha\\in(1,2).\n\\end{equation}\nIt follows from \\cite[Theorem 8]{Ku11} that the Mellin transform $M(s,\\alpha)$ can be written as\n$$\nM(s,\\alpha)=\\alpha^{s-1}\\frac{G(\\alpha\/2;\\alpha)}{G(\\alpha\/2+1;\\alpha)}\\frac{G(\\alpha\/2+2-s;\\alpha)}{G(\\alpha\/2-1+s;\\alpha)}\\frac{G(\\alpha-1+s;\\alpha)}{G(\\alpha+1-s;\\alpha)},\n$$\nwhere $G(z;\\tau)$ is the double gamma function.\n$G(z;\\tau)$ has simple zeroes on the lattice $m\\tau+n$, $m,n\\leq 0$.\nBy a simple calculation one can check that the double gamma functions in the denominators above\nhave no zero with $s=3+iy$ for $|y|\\leq N$.\nIt follows from Lemma \\ref{lemma:double gamma bounded} that there exists a constant $c>0$ such that\n$$\n|M(3+iy,\\alpha)|\\leq ce^{-\\frac{\\pi|y|}{5}} \\text{ for all } |y|\\leq N \\text{ and } \\alpha \\in (1,2),\n$$\nand combining with \\eqref{eqn:Mellin ub} we have \n$$\n|M(3+iy,\\alpha)|\\leq ce^{-\\frac{\\pi|y|}{5}} \\text{ for all } y\\in {\\mathbb R} \\text{ and } \\alpha \\in (1,2).\n$$\n\nHence, \\eqref{eqn:Mellin error} can be estimated as \n$$\n\\left| \\frac{1}{2\\pi i}\\int_{3+i{\\mathbb R}}M(s,\\alpha)x^{-s}ds\\right| \\leq \\frac{x^{-3}}{2\\pi}\\int_{-\\infty}^{\\infty}|M(3+it, \\alpha)|dt\\leq\\frac{x^{-3}}{2\\pi}\\int_{-\\infty}^{\\infty}ce^{-\\frac{\\pi |t|}{5}}dt:=\nAx^{-3}.\n$$\n{\\hfill $\\Box$ \\bigskip}\n\nLet $p^{(\\alpha)}(x)$ be the transition density of $Y^{(\\alpha)}_1$.\nThe following lemma is a variation of \\cite[Propositions 7.1.1 and 7.1.2]{Ko} with an explicit error estimate. \n\\begin{lemma}\\label{lemma:uniform lb}\nLet $\\alpha\\in(1,2)$ and $n\\in{\\mathbb N}$. Then, we have \n$$\np^{(\\alpha)}(x)=\n\\frac{1}{\\pi}\\sum_{k=1}^{n}\\frac{(-1)^{k+1}\\Gamma(1+k\\alpha)\\sin(\\frac{k\\alpha\\pi}{2}) }{k!}x^{-1-\\alpha k}+E(x,\\alpha),\n$$\nwhere \n$$\n\\left|E(x,\\alpha)\\right|\\leq \\frac{2\\Gamma(2n+2)}{\\pi n!}x^{-2-n} \\text{ for all } x\\geq 1 \\text{ and } \\alpha \\in(1,2).\n$$\nIn particular, for any $x\\geq 1$ and $\\alpha\\in (1,2)$ we have\n$$\n \\frac{\\Gamma(1+\\alpha)\\sin\\frac{\\pi\\alpha}{2}}{\\pi}x^{-1-\\alpha}-\\frac{12}{\\pi}x^{-3} \\leq p^{(\\alpha)}(x)\\leq \\frac{\\Gamma(1+\\alpha)\\sin\\frac{\\pi\\alpha}{2}}{\\pi}x^{-1-\\alpha}+\\frac{12}{\\pi}x^{-3},\n$$\nand in fact \n$$\n \\frac{\\Gamma(1+\\alpha)\\sin\\frac{\\pi\\alpha}{2}}{\\pi}=C\\alpha,\n$$\nwhere $C$ is the constant in \\eqref{eqn:same asymptotic}.\n\\end{lemma}\n\\noindent{\\bf Proof.} \nTo prove this lemma, we consider a 1-dimensional $\\alpha$-stable processes $Y^{(\\alpha), \\gamma}_{t}$ with\n$$\n{\\mathbb E}[e^{iy Y_{t}^{(\\alpha), \\gamma}}]=\\exp(-t|y|^{\\alpha}e^{\\frac{i\\pi\\gamma}{2}\\text{sgn}(y)}),\n$$\nwhere $\\gamma$ represents the skewness of $Y^{(\\alpha), \\gamma}_{t}$ and $\\text{sgn}(y):=1_{\\{y\\geq 0\\}}- 1_{\\{y<0\\}}$. \nWe use \n$p^{(\\alpha)}(x, \\gamma)$\nto denote the density of $Y^{(\\alpha), \\gamma}_{1}$. When $\\gamma=0$, $Y^{(\\alpha), \\gamma}_{t}$ reduces to $Y^{(\\alpha)}_{t}$.\nBy the Fourier inversion we have \n\\begin{eqnarray*}\n&&p^{(\\alpha)}(x,\\gamma)\n=\\frac{1}{2\\pi}\\int_{-\\infty}^{\\infty}e^{-ixy}\\exp(-|y|^{\\alpha}e^{\\frac{i\\pi\\gamma}{2}\\text{sgn}(y)})dy\\\\\n&=&\\frac{1}{2\\pi}\\int_{0}^{\\infty}e^{-ixy}\\exp(-|y|^{\\alpha}e^{\\frac{i\\pi\\gamma}{2}})dy+\\frac{1}{2\\pi}\\int_{-\\infty}^{0}e^{-ixy}\\exp(-|y|^{\\alpha}e^{-\\frac{i\\pi\\gamma}{2}})dy\\\\\n&=&\\frac{1}{2\\pi}\\int_{0}^{\\infty}e^{-ixy}\\exp(-|y|^{\\alpha}e^{\\frac{i\\pi\\gamma}{2}})dy+\\frac{1}{2\\pi}\\int^{\\infty}_{0}e^{ixy}\\exp(-|y|^{\\alpha}e^{-\\frac{i\\pi\\gamma}{2}})dy\\\\\n&=&\\frac{1}{\\pi}\\text{Re}\\int_{0}^{\\infty}e^{-ixy}\\exp(-y^{\\alpha}e^{\\frac{i\\pi\\gamma}{2}})dy,\n\\end{eqnarray*} \nwhere $\\text{Re}(\\cdot)$ represents the real part of the argument. \nBy the Taylor expansion with remainder we can write $e^{-ixy}$ as\n$$\ne^{-ixy}=\\sum_{k=0}^{n-1}\\frac{(-ix)^{k}}{k!}y^{k} +\\frac{x^{n}y^{n}}{n!}R,\n$$\nwhere $|R|\\leq 1$.\nSince, \n$$\n\\int_{0}^{\\infty}y^{\\beta}\\exp(-\\lambda y^{\\alpha})dy=\\alpha^{-1}\\lambda^{-\\frac{\\beta+1}{\\alpha}}\\Gamma(\\frac{\\beta+1}{\\alpha}),\n$$\nwe have \n\\begin{align}\\label{eqn:HK}\n&p^{(\\alpha)}(x,\\gamma)\n=\\frac{1}{\\pi}\\text{Re}\\int_{0}^{\\infty}\\left( \\sum_{k=0}^{n-1}\\frac{(-ix)^{k}}{k!}y^{k} +\\frac{x^{n}y^{n}}{n!}R\\right)\\exp(-y^{\\alpha}e^{\\frac{i\\pi\\gamma}{2}})dy\\nonumber\\\\\n&=\\frac{1}{\\pi}\\text{Re}\\sum_{k=0}^{n-1}\\frac{(-ix)^{k}}{k!}\\frac{1}{\\alpha}(e^{\\frac{i\\gamma\\pi}{2}})^{-\\frac{k+1}{\\alpha}}\\Gamma(\\frac{k+1}{\\alpha})+\n\\text{Re}\\left(\\frac{x^{n}}{\\pi\\alpha n!}\\Gamma(\\frac{n+1}{\\alpha})(e^{\\frac{i\\gamma\\pi}{2}})^{-\\frac{n+1}{\\alpha}}R\\right)\\nonumber\\\\\n&=\\frac{1}{\\alpha\\pi}\\text{Re}\\sum_{k=1}^{n}\\frac{(-ix)^{k-1}}{(k-1)!}e^{-\\frac{i\\gamma\\pi k}{2\\alpha}}\\Gamma(\\frac{k}{\\alpha})\n+\\text{Re}\\left(\\frac{x^n}{\\pi\\alpha n!}\\Gamma(\\frac{n+1}{\\alpha})e^{-\\frac{i\\gamma\\pi(n+1)}{2\\alpha}}R\\right)\\nonumber\\\\\n&=\\frac{-1}{\\alpha\\pi x}\\text{Re}\\sum_{k=1}^{n}\\frac{(-x)^{k}}{(k-1)!}\\exp\\left(-i(\\frac{k\\pi(\\gamma-\\alpha)}{2\\alpha} +\\frac{\\pi}{2})\\right)\\Gamma(\\frac{k}{\\alpha})+\\text{Re}\\left(\\frac{x^n}{\\pi\\alpha n!}\\Gamma(\\frac{n+1}{\\alpha})e^{-\\frac{i\\gamma\\pi(n+1)}{2\\alpha}}R\\right)\\nonumber\\\\\n&:=\\frac{1}{\\pi x}\\sum_{k=1}^{n}\\frac{(-x)^{k}}{k!}\\sin(\\frac{k\\pi(\\gamma-\\alpha)}{2\\alpha})\\Gamma(1+\\frac{k}{\\alpha})+R(x,\\alpha),\n\\end{align}\nwhere \n\\begin{equation}\\label{eqn:remainder}\n\\left| R(x,\\alpha)\\right|\\leq \\frac{x^n}{\\pi\\alpha n!}\\Gamma(\\frac{n+1}{\\alpha}).\n\\end{equation}\n\nIt follows from Zolotarev's identity (see \\cite[(7.10)]{Ko}) for $\\alpha\\in(\\frac12,1)\\cup (1,2)$,\n$$\np^{(\\alpha)}(x,\\gamma)=x^{-1-\\alpha}p^{(1\/\\alpha)}(x^{-\\alpha}, \\frac{\\gamma+1}{\\alpha}-1).\n$$\nHence, from \\eqref{eqn:HK} we have \n\\begin{eqnarray*}\n&&p^{(\\alpha)}(x,\\gamma)=x^{-1-\\alpha}p^{(1\/\\alpha)}(x^{-\\alpha}, \\frac{\\gamma+1}{\\alpha}-1)\\\\\n&=&x^{-1-\\alpha}\\left(\\frac{1}{\\pi x^{-\\alpha}}\\sum_{k=1}^{n}\\frac{(-x^{-\\alpha})^{k}}{k!}\\sin(\\frac{k\\pi(\\gamma-\\alpha)}{2})\\Gamma(1+k\\alpha) +R(x^{-\\alpha}, \\frac{1}{\\alpha})\\right).\n\\end{eqnarray*}\nNow let $\\gamma=0$ and define $E(x,\\alpha):=x^{-1-\\alpha}R(x^{-\\alpha}, \\frac{1}{\\alpha})$. From \\eqref{eqn:remainder} we have \n$$\n\\left|E(x,\\alpha)\\right|\\leq \\frac{\\alpha x^{-(n+1)\\alpha-1}}{\\pi n!}\\Gamma(\\alpha(n+1)).\n$$\nHence, for all $\\alpha\\in (1,2)$ and $x\\geq 1$ we have \n$$\n\\left|E(x,\\alpha)\\right|\\leq \\frac{2 x^{-n-2}}{\\pi n!}\\Gamma(2(n+1)).\n$$\n\nThe fact $\\frac{\\Gamma(1+\\alpha)\\sin\\frac{\\pi\\alpha}{2}}{\\pi}=C\\alpha$ can be proved using a similar argument as in the proof of Lemma \\ref{lemma:uniform ub}.\n{\\hfill $\\Box$ \\bigskip}\n\n\n\\begin{lemma}\\label{lemma:Cauchy cancel}\nLet $\\alpha\\in (1,2)$ and suppose that \n$Y^{(\\alpha)}\\in C_{k,k+1}$ for some $k\\in \\mathbb{N}$. \nThere exists a function $\\phi$, independent of $\\alpha\\in (1,2)$, such that \n$$\n{\\mathbb P}(\\overline{Y}^{(\\alpha)}_1>u, Y^{(\\alpha)}_1\\leq u)\\leq \\phi(u) \\text{ with } \\int_{0}^{\\infty}\\phi(u)du<\\infty.\n$$\n\\end{lemma}\n\\noindent{\\bf Proof.} \nNote that \n$$\n{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, Y_{1}^{(\\alpha)}\\leq u)={\\mathbb P}(\\overline{Y}^{(\\alpha)}_1>u)-{\\mathbb P}(Y_{1}^{(\\alpha)}>u).\n$$\nIt follows from Lemmas \\ref{lemma:uniform ub} and \\ref{lemma:uniform lb} that\n$$\n\\overline{p}^{(\\alpha)}(x)\\leq C\\alpha x^{-1-\\alpha} +\nAx^{-3} \\text{ and } \np^{(\\alpha)}(x)\\geq C\\alpha x^{-1-\\alpha} -\\frac{12}{\\pi}x^{-3},\n$$\nwhere $A$ is independent of $\\alpha\\in (1,2)$.\n\nHence, we have \n\\begin{eqnarray*}\n&&{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, Y_{1}^{(\\alpha)}\\leq u)={\\mathbb P}(\\overline{Y}^{(\\alpha)}_1>u)-{\\mathbb P}(Y_{1}^{(\\alpha)}>u)\\\\\n&\\leq&\\int_{u}^{\\infty}\\left(C\\alpha x^{-1-\\alpha} +\nAx^{-3}-C\\alpha x^{-1-\\alpha} +\\frac{12}{\\pi}x^{-3} \\right)dx=\\frac{A}{2}u^{-2}+\\frac{6}{\\pi}u^{-2}.\n\\end{eqnarray*}\nFinally, we set \n$$\n\\phi(u):=1_{\\{01\\}}\n(\\frac{A}{2}+\\frac{6}{\\pi})u^{-2}.\n$$\n{\\hfill $\\Box$ \\bigskip}\n\n\\begin{prop}\\label{prop:Cauchy ub}\nLet $D$ be a bounded $C^{1,1}$ open set in ${\\mathbb R}^{d}$, $d\\geq 2$. Then, we have\n$$\n\\limsup_{t\\to 0}\\frac{|D|-Q_{D}^{(1)}(t)}{t\\ln(1\/t)}\\leq \\frac{|\\partial D|}{\\pi}.\n$$\n\\end{prop}\n\\noindent{\\bf Proof.} \nFrom Lemma \\ref{lemma:Cauchy ub} for any $\\alpha\\in(1,2)$ we have \n\\begin{equation}\\label{eqn:Cauchy ub main}\n\\limsup_{t\\to 0}\\frac{|D|-Q_{D}^{(1)}(t)}{t\\ln(1\/t)}\\leq \\frac{|\\partial D|}{\\pi} +\n\\frac{\\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, Y_{1}^{(\\alpha)}1.\n\\end{equation}\nNote that from \\eqref{eqn:ckl} for \n$Y^{(\\alpha)}\\in C_{k,k+1}$\nwe have $\\alpha \\in (1,2)$ and as $k\\to \\infty$, $\\alpha\\downarrow 1$.\nHence, it follows from Lemma \\ref{lemma:Cauchy cancel} and the Lebesgue dominated convergence theorem that\n$$\n\\lim_{Y_{t}^{(\\alpha)}\\in C_{k,k+1}, k\\to \\infty}\\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, Y_{1}^{(\\alpha)}\\leq u)du\n=\\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(1)}_{1}>u, Y_{1}^{(1)}\\leq u)du.\n$$\nNote that the density $p^{(1)}(x)$ of $Y^{(1)}_{1}$ is given by \n$$\np^{(1)}(x)=\\frac{1}{\\pi(1+x^2)}\n$$\nand $\\int_{u}^{\\infty}{\\mathbb P}(Y_{1}^{(1)}>u)du=\\int_{u}^{\\infty}\\frac{1}{\\pi (1+x^2)}dx=\\frac{1}{\\pi}(\\frac{\\pi}{2}-\\arctan(u))=\\frac{1}{\\pi}\\arctan(\\frac{1}{u})$.\nBy the Taylor expansion of $\\arctan x =\\sum_{n=1}^{\\infty}(-1)^{n+1}\\frac{x^{2n-1}}{2n-1}$ for $|x|<1$ we have $\\frac{1}{\\pi}\\arctan(1\/u)\\geq \\frac{1}{\\pi u} -\\frac{1}{2\\pi u^3}$ for all sufficiently large $u$.\nHence, it follows from \\cite[(3.5)]{P20} that \n\\begin{eqnarray*}\n&&{\\mathbb P}(\\overline{Y}^{(1)}_{1}>u, Y_{1}^{(1)}\\leq u ) ={\\mathbb P}(\\overline{Y}^{(1)}_{1}>u)-{\\mathbb P}(Y_{1}^{(1)}> u)\\\\\n&=&{\\mathbb P}(\\overline{Y}^{(1)}_{1}>u)-\\frac{\\arctan(1\/u)}{\\pi}\\leq{\\mathbb P}(\\overline{Y}^{(1)}_{1}>u)-\\frac{1}{\\pi u} +\\frac{1}{2\\pi u^3}\\leq \\frac{4}{\\pi^2}\\frac{\\ln u}{u^3} +\\frac{1}{2\\pi u^3}\n\\end{eqnarray*}\nfor all sufficiently large $u>0$, and this implies \n\\begin{equation}\\label{eqn:Cauchy finite}\n\\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(1)}_{1}>u, Y_{1}^{(1)}\\leq u)du<\\infty.\n\\end{equation}\n\nHence, if follows from \\eqref{eqn:Gamma} and \\eqref{eqn:Cauchy finite} that\n$$\n\\lim_{Y^{(\\alpha)}\\in C_{k,k+1}, k\\to \\infty}\n\\frac{\\int_{0}^{\\infty}{\\mathbb P}(\\overline{Y}^{(\\alpha)}_{1}>u, Y_{1}^{(\\alpha)}\\tilde{T}$ and possibly also two\nadditional dynamical variables---which are characterized by the following\ntwo properties:\n\n(i) For the same variables of the original dynamical system $D$ the new\ndynamical system $\\tilde{D}$ yields, over the time interval $\\tilde{T}$,\nhence over an \\textit{arbitrarily long} time, a dynamical evolution which\nmimics \\textit{arbitrarily closely} that yielded by the original system $D$;\nup to corrections of order $t\/\\tilde{T}$, or possibly even \\textit\nidentically}.\n\n(ii) The system $\\tilde{D}$ is \\textit{isochronous}: for arbitrary initial\ndata all its solutions are completely periodic with a period $T$, which can\nalso be arbitrarily assigned, except for the\\ (obviously necessary)\ncondition $T>\\tilde{T}$.\n\nMoreover it has been shown \\cite{CL2007,calogero} that, if the dynamical\nsystem $D$ is a many-body problem characterized by a (standard, autonomous)\nHamiltonian $H$ which is translation-invariant (i. e., it features no\nexternal forces), other (also autonomous) Hamiltonians $\\tilde{H}$\ncharacterizing \\textit{modified} many-body problems can be manufactured\nwhich feature the \\textit{same} dynamical variables as $H$ (i. e., in this\ncase there is no need to introduce two additional dynamical variables) and\nwhich yield a time evolution \\textit{quite close}, or even \\textit{identical\n, to that yielded by the original Hamiltonian $H$ over the \\textit\narbitrarily assigned} time $\\tilde{T}$, while being \\textit{isochronous}\nwith the \\textit{arbitrarily assigned} period $T$, of course with $T>\\tilde{\n}$.\n\nLet us emphasize that the class of Hamiltonians $H$ for which this result is\nvalid is quite general. In particular it includes the standard Hamiltonian\nsystem describing an \\textit{arbitrary} number $N$ of point particles with\n\\textit{arbitrary} masses moving in a space of \\textit{arbitrary} dimensions\n$d$ and interacting among themselves via potentials depending \\textit\narbitrarily} from the interparticle distances (including the possibility of\nmultiparticle forces), being therefore generally valid for any realistic\nmany-body problem, hence encompassing most of nonrelativistic physics. This\nresult is moreover true, \\textit{mutatis mutandis}, in a quantal context.\n\nFor instance let us tersely review here the case of the standard Hamiltonian\ndescribing the many-body problem, but focusing---merely for notational\nsimplicity---on the case with equal particles (and setting their mass to\nunity) and a one-dimensional setting, which reads (in self-evident standard\nnotation)\n\\begin{subequations}\n\\begin{equation}\nH\\left( \\underline{p},\\underline{q}\\right) =\\frac{1}{2}\\sum_{n=1}^{N}\\left(\np_{n}^{2}\\right) +V\\left( \\underline{q}\\right) \\label{H}\n\\end{equation\nwith the potential being translation-invariant, $V\\left( \\underline{q\n+a\\right) =V\\left( \\underline{q}\\right) $ (i. e., no external forces) but\notherwise unrestricted (except for the condition---again, for\nsimplicity---that the time-evolution yielded by this Hamiltonian be\nnonsingular). The standard approach to treat this problem is to introduce\nthe center-of-mass coordinate $Q\\left( \\underline{q}\\right) =\\frac{1}{N\n\\sum_{n=1}^{N}q_{n}$ and the total momentum $P\\left( \\underline{p}\\right)\n=\\sum_{n=1}^{N}p_{n},$ and to then focus on the coordinate and momenta in\nthe center of mass, $x_{n}=q_{n}-Q\\left( \\underline{q}\\right) ,$ \ny_{n}=p_{n}-P\\left( \\underline{p}\\right) \/N$, which characterize the physics\nof the problem and whose evolution is determined by the Hamiltonian $h\\left(\n\\underline{y},\\underline{x}\\right) $ defined as follows\n\\begin{equation}\nH\\left( \\underline{p},\\underline{q}\\right) =\\frac{\\left[ P\\left( \\underline{\n}\\right) \\right] ^{2}}{2N}+h\\left( \\underline{y},\\underline{x}\\right) ~,\n\\label{Hh}\n\\end{equation\n\\begin{equation}\nh\\left( \\underline{y},\\underline{x}\\right) =\\frac{1}{2}\\sum_{n=1}^{N}\\left(\ny_{n}^{2}\\right) +V\\left( \\underline{x}\\right) ~, \\label{h}\n\\end{equation\nwhere $V\\left( \\underline{x}\\right) =V\\left( \\underline{q}\\right) $ thanks\nto the assumed translation-invariance of $V\\left( \\underline{q}\\right) $. An\n\\textit{isochronous} Hamiltonian $\\tilde{H}\\left( \\underline{p},\\underline{q\n;T\\right) $ reads then as follows,\n\\end{subequations}\n\\begin{subequations}\n\\begin{equation}\n\\tilde{H}\\left( \\underline{p},\\underline{q};T\\right) =\\frac{1}{2\n\\sum_{n=1}^{N}\\left\\{ \\left[ P\\left( \\underline{p}\\right) +h\\left(\n\\underline{y},\\underline{x}\\right) \\right] ^{2}+\\left( \\frac{2\\pi }{T\n\\right) \\left[ Q\\left( \\underline{q}\\right) \\right] ^{2}\\right\\} ~.\n\\label{Htilde}\n\\end{equation\nIt can indeed be shown \\cite{CL2007,calogero} that it entails an \\textit\nisochronous} evolution (with period $T$) of the center of mass coordinate $Q\n, of the total momentum $P$, and---most importantly---of all the particle\ncoordinates, whose evolution then read\n\\begin{equation}\n\\underline{\\tilde{x}}\\left( t\\right) =\\underline{x}\\left( \\tau \\left(\nt\\right) \\right) ~,~~~\\underline{\\tilde{y}}\\left( t\\right) =\\underline{y\n\\left( \\tau \\left( t\\right) \\right) ~, \\label{xytildet}\n\\end{equation\nwhere we indicate as $\\underline{\\tilde{x}}\\left( t\\right) ,$ $\\underline\n\\tilde{y}}\\left( t\\right) $ the time evolution yielded by the tilded\nHamiltonian $\\tilde{H}\\left( \\underline{p},\\underline{q};T\\right) $ and as \n\\underline{x}\\left( t\\right) ,$ $\\underline{y}\\left( t\\right) $ the time\nevolution yielded by the original Hamiltonian $H\\left( \\underline{p}\n\\underline{q}\\right) ,$ an\n\\begin{equation}\n\\tau \\left( t\\right) =A~\\sin \\left( \\frac{2\\pi t}{T}+\\Phi \\right)\n\\label{taut}\n\\end{equation\nwhere $A$ and $\\Phi $ are \\textit{constant} parameters given by simple\nexpressions in terms of the initial values, $Q\\left( 0\\right) $ and $P\\left(\n0\\right) $, of the position of the center of mass of the system and of its\ntotal momentum, and of the Hamiltonian $\\tilde{H}\\left( \\underline{p}\n\\underline{q};T\\right) $ (which is of course a constant of motion for this\ntime evolution). Note that, for $\\left\\vert t\\right\\vert <\\tilde{T}$ in the\ntime coordinate $t\\,$); but it is \\textit{degenerate} at an infinite,\ndiscrete sequence of times $t_{n}=t_{0}\\pm nT\/2,$ $n=0,1,2,...$ . We\ninterpreted these metrics as corresponding to \\textit{isochronous\ncosmologies }\\cite{isochronous cosmologies}.\n\nLet us also recall \\cite{isochronous cosmologies} that this result is not\nrestricted to homogeneous, isotropic and spatially flat metrics: it can be\neasily extended to any synchronous metric, therefore it is quite general\nsince most metrics can be written in synchronous form by a diffeomorphic\nchange of coordinates.\n\nDue to the diffeomorphic correspondence between $g_{\\mu \\nu }$ and $\\tilde{g\n_{\\mu \\nu }$ \\textit{locally} in time---for time intervals of order $\\tilde{\n}$---these two metrics give the same physics locally in time, and therefore\nthere is no way to distinguish them using observations local in time; which\nis essentially the same finding valid in the context of the Hamiltonian\nsystems considered in \\cite{CL2007,calogero} and tersely recalled above. In\nparticular, the metric $\\tilde{g}_{\\mu \\nu }$---while being cyclic on time\nscales larger than $\\tilde{T}$---may yield over the time interval $\\tilde{T}$\njust the standard sequence of domination firstly by radiation, then by\nmatter and dark energy, as well as the cosmological perturbations consistent\nwith all observational tests \\cit\n{cmb,leansing,bao,lss,supernovae,ref0,ref1,ref2,ref3,ref4,bicep2}, which\ncharacterize the metric $g_{\\mu \\nu }$ of the standard $\\Lambda $-CDM\ncosmological model (see for instance \\cite{mukhanov} for a review).\nFurthermore, it has been shown that such isochronous metrics $\\tilde{g}_{\\mu\n\\nu }$ can be manufactured to be geodesically complete \\cite{isochronous\ncosmologies} and therefore singularity-free\\footnote\nA spacetime is singularity-free if it is geodesically complete, i.e. if its\ngeodesics can be always past- and future-extended \\cite{poisson,MTW}.}, so\nthat the geodesic motion as well as all physical quantities described by\nscalar invariants are always well defined \\cite{isochronous cosmologies} and\nthe Big Bang singularity may be avoided---even when some of the\nphenomenological observations can nevertheless be interpreted as remnants of\na past Big Bang, which however may never be actually attained by the metric \n\\tilde{g}_{\\mu \\nu }$, neither in the past nor in the future. This implies\nthat these isochronous metrics $\\tilde{g}_{\\mu \\nu }$ may describe a\nsingularity-free universe which---while featuring a time evolution which\nreproduces identically (up to diffeomorphic time reparameterization; \\textit\nlocally} in time, except at a discrete set of instants $t_{n}$) the standard\ncosmological model characterized by the metric $g_{\\mu \\nu }$---features an\nexpansion which stops at some instant, to be followed by a period of\ncontraction, until this phase of evolution stops and is again followed by an\nexpanding phase, this pattern being repeated \\textit{ad infinitum}. Let us\nre-emphasize that the universe characterized by the metric $\\tilde{g}_{\\mu\n\\nu }$ may thereby avoid to experience the Big Bang singularity, even if the\nuniverse which it mimics locally in time, characterized by the metric \ng_{\\mu \\nu },$ does encounter that singularity. Let us also note that this\nclass of metrics realize \\textit{de facto} the reversal of time's arrow, as\ndiscussed for instance in \\cite{sakharov}.\n\nBut in order to do so, these isochronous metrics must be \\textit{degenerate}\nat a discrete set of instants $t_{n}$ when the time reversals occur, say\nfrom expansion to contraction and viceversa; a phenomenology whose inclusion\nin general relativity might be considered problematic, because at these\ninstants $t_{n}$---when the metric is \\textit{degenerate}---the \"equivalence\nprinciple\" corresponding to the requirement that the metric tensor have a\nMinkowskian signature is violated. On the other hand a \\textit{physically\nunobservable} violation of a \"principle\" can be hardly considered \\textit\nphysically relevant}.\n\nPurpose and scope of this paper is to better clarify the meaning and the\ninterpretation of the \\textit{isochronous spacetimes} introduced in \\cit\n{isochronous cosmologies}, and to discuss some additional technical aspects.\nAs stated above, these isochronous metrics have the property to be \\textit\ndegenerate} at an infinite set of times $t_{n}$, i.e. over an infinite\nnumber of $3$-dimensional hypersurfaces, where generally the scale factor\nreaches its maximum and minimum values and the spacetime changes from an\nexpanding to a contracting phase and vice versa \\cite{isochronous\ncosmologies}.\n\nThe inclusion of degenerate metrics in general relativity is not a trivial\nmatter; indeed, it is forbidden if it were required---as an absolute (as it\nwere, \"metaphysical\") rule---that general relativity be consistent with the\n\"equivalence principle\", i. e. that spacetime \\textit{always} be a\nPseudo-Riemannian manifold \\textit{with Minkowskian signature}. This would\nindeed imply that the solutions of Einstein's general relativity, to be\nacceptable, must be \\textit{nondegenerate} metric tensors, since wherever in\nspacetime a metric is \\textit{degenerate} it is not possible to define its\ninverse, hence the Christoffel symbols as well as the Riemann, Ricci and\nEinstein tensors are not defined there. Hence---so the argument of some\ncritics goes---our isochronous metrics, which are not Minkowskian (locally,\non the hypersurfaces $t=t_{n}$) should be discarded in general relativity\nbecause they violate the equivalence principle at those degeneracy surfaces.\nOur rejoinder is that the metrics we introduced \\cite{isochronous\ncosmologies} are solutions of a version of general relativity whose\ndynamics, while still governed by Einstein's equations\\footnote\nEinstein's equations by itself do not determine the signature of spacetime,\nsee the discussion in \\cite{teitelboim,ellis1}.}, does allow the violation\nof the equivalence principle at a discrete set of hypersurfaces. It can be\nargued that such theories are therefore \\textit{different} from standard\ngeneral relativity, if validity of the equivalence principle---enforced as\nan absolute, universal rule---is considered an essential element of that\ntheory; although the difference is in fact \\textit{physically unobservable}.\nThis situation is indeed quite analogous with the findings described above,\nvalid in the nonrelativistic context of Hamiltonian many-body problems\n(classical or quantal), where the solutions of a physical theory described\nby a Hamiltonian $H$ are arbitrarily well approximated or even identically\nreproduced over an arbitrary time interval by periodic solutions of a\n\\textit{different} Hamiltonian $\\tilde{H}$, i.e. by a \\textit{different}\nphysical theory. The main difference is that in the general relativistic\ncontext in\\textbf{\\ }the two \\textit{different} physical theories the\ndynamic is given by the\\textit{\\ same} Einstein's equations, only the class\nof acceptable solutions is \\textit{different}.\n\nThe consideration in the context of general relativity of degenerate metrics\nis in any case not a novelty. For instance they were already introduced \\cit\n{ellis1,ellis2}, and subsequently investigated in a series of papers \\cit\n{elliscarfora,elliscomment,ellis3,ellis4,ellis5,ellis6,ellis8,ellis9,ellis10,ellis11,ellis12,ellis13,ellis14,ellis15,ellis16,ellis17,ellis18,ellis19}\n, in order to investigate signature-changing spacetimes. This class of\nsignature-changing metrics corresponds to a classical realization of the\nchange of signature in quantum cosmology conjectured by Hartle and Hawking\n\\cite{hartlehawking3} (see also \\cit\n{hartlehawking1,hartlehawking2,hartlehawking4}), which has philosophical\nimplications on the origin of the universe \\cite{hawkinghistory}.\n\nBelow we will also consider a different realization of isochronous\ncosmologies via non-degenerate metrics featuring a jump in their first\nderivatives at the inversion times $t_{n}$, which then implies a\ndistributional contribution in the stress-energy tensor at $t_{n}$. These\nmetrics are diffeomorphic to the continuous and degenerate isochronous\nmetrics everywhere except on the inversion hypersurfaces $t=t_{n}$, hence\nthey describe the same physics except at the infinite set of discrete times \nt_{n};$ therefore they may also be manufactured so as to agree with all\ncosmological observations (\\textit{locally} in time). Since these two\nrealizations of isochronous cosmologies are \\textit{locally} but \\textit{not\nglobally} diffeomorphic, they correspond to different spacetimes and may be\nconsidered to emerge from \\textit{different} theories: the nondegenerate\nones are generalized (in the sense of distributions) solutions of\nEinsteinian general relativity (including universal validity of the\nequivalence principle, but allowing jumps---whose physical significance and\njustification is moot---at the discrete times $t_{n}$); the degenerate ones\nfeature no mathematical or physical pathologies but fail to satisfy the\n\"equivalence principle\" at the discrete times $t_{n}$. They are essentially\nindistinguishable---among themselves and from non-isochronous\ncosmologies---by feasible experiments (unless one considers feasible an\nexperiment lasting an arbitrarily long time...).\n\n\\section{Isochronous Cosmologies}\n\nThe procedure used to construct, for an autonomous dynamical system $D$,\nanother autonomous dynamical system $\\tilde{D}$ whose solutions \\textit\napproximate} arbitrarily closely, or even reproduce exactly, those of the\nsystem $D$ over an \\textit{arbitrary} time interval $\\tilde{T}$ but are\n\\textit{isochronous} with an \\textit{arbitrary} period $T>\\tilde{T}$, is\nbased on the introduction of an auxiliary variable $\\tau $ in place of the\nphysical time variable $t$ \\cite{CL2007,calogero}. This procedure formally\nentails a change in the time dependence of any physical quantity $f(t)$, so\nthat $f(t)$ gets replaced by $f(\\tau \\left( t\\right) )$ with $\\tau \\left(\nt\\right) $ a periodic function of $t$ (with period $T$: $\\tau \\left( t\\pm\nT\\right) =\\tau \\left( t\\right) $), hence the new dynamics entails that all\nphysical quantities evolve periodically with period $T$. This change, in the\ncase of a nonrelativistic dynamical system, implies a change in the\ndynamical equations: for instance, in the case of the (quite general)\nmany-body problem described above, the theory is characterized by a new\n(autonomous) Hamiltonian $\\tilde{H}$---different from the original\n(autonomous) Hamiltonian $H$---which causes any physical variable originally\nevolving as $f\\left( t\\right) $ under the dynamics yielded by $H$, to evolve\ninstead as $f(\\tau (t))$ under the dynamics yielded by $\\tilde{H}$. This\nchange of the physical laws determining the time evolution of the system\nproduces the \\textit{isochronous} evolution \\cite{calogero,CL2007}. But let\nus stress---as we already did in \\cite{isochronous cosmologies}---that any\nattempt to attribute to the new variable $\\tau $ the significance of \"time\"\nwould be improper and confusing: the variable playing the role of \"time\" is\nthe same, $t$, for both the original dynamical system $D$ and the modified\ndynamical system $\\tilde{D}$; in particular, for that characterized by the\nstandard many-body Hamiltonian $H,$ see (\\ref{H}), as well as for that\ncharacterized by the modified many-body Hamiltonian $\\tilde{H}$, see (\\re\n{Htilde}).\n\nThe dynamical systems $D$ and $\\tilde{D}$ mentioned above feature the same\ntrajectories in phase space, but while the time evolution of the dynamical\nsystem $D$ corresponds to a \\textit{uniform} forward motion along those\ntrajectories, the time evolutions of the modified dynamical systems $\\tilde{\n}$---although produced by time-independent equations of motion---correspond\nto a \\textit{periodic} (with assigned period $T$), forward and backward,\ntime evolution along those same trajectories, exploring therefore only a\nportion of them; with the possibility to manufacture $\\tilde{D}$ in such a\nway that, for an \\textit{arbitrary} subinterval $\\tilde{T}0$ for any $t\n. This choice implies that all the scalar invariants as well as the pressure\nand energy density of the fluid remain finite at any $t$ and the spacetime\nis geodesically complete, thus singularity free.\n\n(vi) Geodesics have, in spacetime, a helical structure.\n\nThese solutions of course violate the \"equivalence principle\" at the\ninfinite set of discrete times $t_{n}.$ It is thus seen that, with the\nchoice (\\ref{bBperiodic}), the class of spacetime solutions of the equations\nof general relativity is enlarged by allowing unobservable violations of the\n\"equivalence principle\" at a discrete, numerable set of times $t_{n}$. This\nmore general class of spacetimes do not seem to entail any \\textit\nexperimentally observable} differences; although they might entail a quite\ndifferent structure of spacetime, for instance absence versus presence of\nBing Bang singularities.\n\nThe problem of finding a generalization of general relativity which includes\ndegenerate metrics has been studied in the past, motivated by the\nconsideration of a class of metrics, the so called signature-changing\nmetrics \\cit\n{ellis1,ellis2,elliscarfora,elliscomment,ellis3,ellis4,ellis5,ellis6,ellis8,ellis9,ellis10,ellis11,ellis12,ellis13,ellis14,ellis15,ellis16,ellis17,ellis18,ellis19}\n, which give a classical realization of the change of signature in quantum\ncosmology conjectured by Hartle and Hawking \\cit\n{hartlehawking3,hartlehawking1,hartlehawking2,hartlehawking4,hawkinghistory\n. The classical change of signature for a homogeneous, isotropic and\nspatially flat universe is realized by a metric tensor defined by the\nfollowing line elemen\n\\begin{equation}\nds^{2}=N(t)~dt^{2}-R(t)^{2}~d\\vec{x}^{2}~,\n\\label{changeofsignaturedegenerate}\n\\end{equation\nwhere $N(t)$ is a continuous function changing sign at some time $t_{0}$,\nfor instance $N(t)$ is positive for $t>t_{0}$, negative for $t\\tau _{0}$ and \nf(\\tau )=-1$ for $\\tau <\\tau _{0}$. It has been shown \\cite{ellis1,ellis2}\nthat in this case it is possible to introduce a smooth (generalized)\northonormal reference frame which allows a variational derivation of\nEinstein's equations and an extension of the Darmois formalism \\cite{ellis16}\nto discontinuous metrics, in such a way that, also in the case of a\ndiscontinuous metric tensor, the discontinuity of the extrinsic curvature on\na hypersurface is related to a distributional stress energy tensor \\cit\n{ellis15}. This is achieved by adding a surface term to the Einstein-Hilbert\naction, in such a way that the gravitational action become\n\\begin{equation}\nS_{g}=\\int_{M\/\\Sigma }\\sqrt{-g}Rd^{4}x+\\oint_{\\Sigma }\\sqrt{-g}Kd\\Sigma\n\\label{actionboundary}\n\\end{equation\nwhere $M$ is the manifold defining the spacetime, $\\Sigma $ is the boundary\nof $M$ where the metric tensor is discontinuous, $R$ is the Ricci scalar\ncurvature and $K$ the extrinsic curvature of $\\Sigma $. The addition of such\na boundary term is always necessary when one considers manifolds with\nboundaries \\cite{HawkingHorowitz}. Incidentally we note that this action\ncannot be used in the case of degenerate metrics, because in this case the\ndegeneracy surface $\\Sigma $ has no unitary normal vector hence the\nextrinsic curvature $K$ does not exist.\n\nThe two realizations of the classical change of signature, the continuous\nand degenerate one, see (\\ref{changeofsignaturedegenerate}), and the\ndiscontinuous one, see (\\ref{changeofsignaturediscontinuous}), are \\textit\nlocally but not globally} diffeomorphic, since they are related by a change\nof the time variable $dt=d\\tau \/\\sqrt{|N(t)|}$ which is not defined at the\ntime of signature change $t_{0}$ when $N(t)=0$. Hence they represent two\ndifferent spacetimes. However, except for the instants when $N(t)$ vanishes,\nthey are \\textit{locally } diffeomorphic and thus they describe---excepts at\nthose instants---the same physics .\n\nAs in the case of signature-changing metrics, it is also possible to\nconsider a realization of \\textit{isochronous} cosmologies different from\nthat we introduced in \\cite{isochronous cosmologies} and reviewed above (see\n(\\ref{newmetric}) and the subsequent treatment), via \\textit{nondegenerate}\nmetrics featuring a finite jump of their first derivatives at an infinite,\ndiscrete set of equispaced times. Such a realization is given by the\nfollowing metric\n\\begin{equation}\nds^{2}=d\\eta ^{2}-\\tilde{a}(\\eta )^{2}d\\vec{x}^{2}\n\\label{newmetricnondegenerate}\n\\end{equation\nwith\n\\begin{equation}\n\\alpha (\\tau (\\eta ))\\equiv \\tilde{a}(\\eta )\n\\label{newscalarfactornondegenerate}\n\\end{equation\nand\n\\begin{equation}\nd\\tau =C(\\eta )~d\\eta , \\label{newCndegenerate}\n\\end{equation\nwith a periodic function $C(\\eta )$ of period $T$ given for instance by \nC(\\eta )=1$ for $nT<\\eta <(2n+1)T\/2$ and $C(\\eta )=-1$ for $(2n+1)T\/2<\\eta\n<(n+1)T$ and $n$ integer, hence such that\n\\begin{subequations}\n\\begin{equation}\n\\tau \\left( \\eta \\right) \\equiv \\int_{0}^{\\eta }C\\left( \\eta ^{\\prime\n}\\right) d\\eta ^{\\prime }=\\eta -nT~~~\\text{for}~~~nT0\\, ,\\quad 3\\lambda + 2\\mu > 0\\, ,\\quad \\eta > 0\\, .\\]\nOn the other hand, the condition \\(u_i \\neq 0\\), requires \\(\\mathcal{\\tilde T}\\)\nto be different from zero and thus the condition required by\n\\eqref{eq:definiteness}. In the case of rigid body motion, \\(\\mathcal{U}\\)\ncould be zero implying that the form is positive semi-definite, while the form\n\\(\\mathcal{\\tilde T}\\) is positive definite.\n\n\\section{Finite element formulation}\nIn this section, we derive a consistent finite element formulation for periodic\ncouple stress elastodynamics, as an extension of those formulated by\n\\cite{darrall2014} for the corresponding quasistatic problem and by\n\\cite{guarin2020} for periodic micropolar Bloch analysis. In particular,\nthe \\(C^1\\) displacement continuity requirement is avoided by using a Lagrange\nmultiplier approach. Other finite element solutions in C-CST include a penalty\nmethod for isotropic elastostatics \\citep{chakravarty2017}, Lagrange multipliers\nfor centrosymmetric anisotropic elastostatics \\citep{pedgaonkar2021} and mixed\nvariable methods for isotropic elastodynamics \\citep{deng2016, deng2017}.\n\n\\subsection{Lagrange multiplier reformulation}\nConsider now a modification of the action given in\n\\eqref{eq:hamilton_freq_free} to include Lagrange multipliers \\(\\lambda_i\\) that\nenforce compatibility between the displacement field \\(u_i\\) and an assumed\nindependent rotation field \\(\\theta_i\\). Thus, the modified action becomes\n\\begin{equation}\n\\begin{split}\n\\mathcal{\\hat A}[u; \\omega] = &\\frac{1}{2}\\int\\limits_{V} e_{ij}^* C_{ijkl} e_{ij} \\dd{V}\n + \\frac{1}{2}\\int\\limits_{V} \\kappa_{i}^* D_{ij} \\kappa_{i } \\dd{V}\n -\\frac{\\omega^2}{2}\\int\\limits_{V} u_i^* \\rho u_i \\dd{V}\\\\\n&- \\int\\limits_{S_t} u_i^* t_i \\dd{S} - \\int\\limits_{S_m} \\theta_i^* m_i \\dd{S}\\\\\n&+ \\int\\limits_{V}\\lambda_i^* (\\epsilon_{ijk} u_{k,j} - 2\\theta_i )\\dd{V}\\, .\n\\end{split}\n\\label{eq:hamilton_freq_with_multipliers}\n\\end{equation}\n\nFor stationarity, we require\n\\[\\var{\\hat{\\mathcal{A}}} = \\pdv{\\hat{\\mathcal{A}}}{u_i} \\var{u_i}\n+ \\pdv{\\hat{\\mathcal{A}}}{\\theta_i}\\var{\\theta_i}\n+ \\pdv{\\hat{\\mathcal{A}}}{\\lambda_i} \\var{\\lambda_i} = 0\\, , \\]\nwhich is equivalent to\n\\begin{equation}\n\\begin{split}\n&\\int\\limits_{V} \\delta e_{ij}^* C_{ijkl} e_{ij} \\dd{V}\n + \\int\\limits_{V} \\delta \\kappa_{i}^* D_{ij} \\kappa_{i } \\dd{V}\n -\\omega^2 \\int\\limits_{V} \\delta u_i^* \\rho u_i \\dd{V}\\\\\n&\\qquad- \\int\\limits_{S_t} \\delta u_i^* t_i \\dd{S} - \\int\\limits_{S_m} \\delta \\theta_i^* m_i \\dd{S}\\\\\n&\\qquad+ \\int\\limits_{V}\\delta \\lambda_i^* (\\epsilon_{ijk} u_{k,j} - 2\\theta_i )\\dd{V}\\\\\n&\\qquad+ \\int\\limits_{V} (\\epsilon_{ijk} \\delta u_{k,j}^* - 2\\delta \\theta_i^* ) \\lambda_i \\dd{V}\\, .\n\\end{split}\n\\label{eq:weak_form_with_multipliers}\n\\end{equation}\n\n\\Cref{eq:weak_form_with_multipliers} is the modified weak form that will be used\nhere as the basis for the finite element Bloch analysis of an elastic\ncouple-stress solid. The Lagrange multiplier terms enforce the required\nkinematic constraint between the continuum rotations \\(\\epsilon_{ijk} u_{k,j}\\)\nof the material point and the independent rotational variables \\(\\theta_i\\).\n\nFrom \\eqref{eq:weak_form_with_multipliers}, we obtain the following\nEuler-Lagrange equations\n\\begin{equation}\n\\begin{split}\n&(C_{ijkl} e_{kl} + \\epsilon_{ijk} \\lambda_k)_{,j} = -\\rho \n\\omega^2u_i\\quad \\text{in } V,\\\\\n&\\epsilon_{ijk} (D_{kl} \\kappa_{l})_{,j} - 2\\lambda_i = 0\\quad \\text{in } V,\\\\\n&\\theta_i = \\frac{1}{2}\\epsilon_{ijk} u_{k,j}\\quad \\text{in } V,\\\\\n&t_i = (C_{ijkl} e_{kl} + \\epsilon_{ijk}\\lambda_k)n_j\\quad \\text{on } S_t,\\\\\n&m_i = \\epsilon_{ijk} D_{kl} \\kappa_{l} n_j\\quad \\text{on } S_m\\, ,\n\\end{split}\n\\label{eq:bvp}\n\\end{equation}\nComparing this with \\eqref{eq:skew_symmetry}, we can conclude that the Lagrange\nmultipliers equal the skew-symmetric part of the force-stress tensor, i.e.,\n\\[\\lambda_i = s_i\\, .\\]\n\n\\subsection{Discretization}\nTo discretize \\eqref{eq:weak_form_with_multipliers}, we use for the\nelement-based shape functions second-order Lagrange interpolation for the\ndisplacements and rotations and constant skew-symmetric stresses. This\ntranslates into \\(C^0\\) inter-element displacement and rotation continuity, and\nskew-symmetric stresses that are constant within the element but discontinuous\nbetween elements. \\Cref{fig:element} depicts a typical element for the\ndiscretization and the degrees of freedom used in two-dimensional idealizations.\n\\begin{figure}[H]\n\\centering\n\\includegraphics[width=2 in]{img\/element.pdf}\n\\caption{Finite element used for the finite element discretization of the\nC-CST material model. A second-order Lagrange interpolation is used for\ndisplacements and rotations while a constant is used for the skew-symmetric\nstresses. The nodal degrees of freedom are depicted as black disks while the\nwhite disk describes the element skew-symmetric force-stress degree of freedom.}\n\\label{fig:element}\n\\end{figure}\n\nTo write the discretized equations, we will use a combined index notation. In\nthis context subscripts will still make reference to scalar components of\ntensors while capital superscripts will indicate interpolation operations. For\ninstance in the expression\n\\[u_i = _{u}\\hspace{-4pt}N_{i}^{Q}u^{Q}\\]\nsubscripts $i$ indicate the scalar components of the vector \\(u\\). To facilitate\nfurther operations this subscript is also placed in the shape function resulting\nin terms like \\(_{u}N_{i}^{Q}\\) and where the term \\(u^{Q}\\) represents the\nnodal point displacement associated to the $Q$th nodal point. This nodal vector\nimplicitly considers horizontal and vertical rectangular components. To clarify,\nthe displacement interpolation scheme written here as\n\\(u_i = _{u}\\hspace{-4pt}N_{i}^{Q}u^{Q}\\) takes the following explicit form for\nthe single nodal point Q:\n\\begin{equation}\n\\begin{bmatrix}u_x\\\\ u_y\\end{bmatrix}=\n\\left[\\cdots\\begin{array}{cc}N^Q &0\\\\ 0 &N^Q\\end{array}\\cdots\\right]\n\\begin{Bmatrix}\\vdots\\\\ u_x^Q\\\\ u_y^Q\\\\\\vdots\\end{Bmatrix} \\, .\n\\label{eq:interp_nota}\n\\end{equation}\n\nWith this notation we write for the primary variables\n\\((u_i, \\theta_i, s_i )\\) the following interpolated versions\n\\begin{equation}\nu_i = _{u}\\hspace{-4pt}N_{i}^{Q}u^{Q},\\quad \\theta_i\n = _{\\theta}\\hspace{-4pt}N_i^{Q}\\theta^{Q} ,\\quad s_i\n = _{s}\\hspace{-4pt}N_{i}^{Q} s^{Q} \\, ,\n\\label{eq:interp_fun_prim}\n\\end{equation}\nand similarly for the secondary kinematic descriptors\n\\(e_{ij} , \\epsilon_{ijk} u_{i,j}\\) and \\(\\kappa_{i}\\) \n\\begin{equation}\ne_{ij} = _{e}\\hspace{-4pt}B_{ij}^{Q}u^{Q} ,\\quad \\epsilon_{ijk} u_{i,j}\n = _{\\nabla}\\hspace{-4pt}B_{k}^{Q}u^{Q} ,\\quad \\kappa_{i}\n = _{\\kappa}\\hspace{-4pt}B_{i}^Q \\theta^Q \\, ,\n\\label{eq:interp_fun_second}\n\\end{equation}\ntogether with the constitutive equations\n\\begin{equation}\n\\begin{aligned}\n &\\sigma_{ij} = C_{ijkl}e_{kl}\\, , \\\\\n &\\mu_{i} = D_{ij}\\kappa_{j}\\, .\n\\end{aligned}\n\\label{eq:const}\n\\end{equation}\n\n\nSubstitution of the above relations in \\eqref{eq:weak_form_with_multipliers}\ngives the discrete version of the first variation of the modified correlated\naction;\n\\begin{equation}\n\\begin{split}\n&\\var{\\hat{\\mathcal{A}}} = \\delta u^{Q*} \\int\\limits_{V} (_{e}B_{ij}^{Q}) (C_{ijkl}) (_{e}B_{kl}^{P}) \\dd{V} u^{P} - \\rho \\omega ^2 \\delta u^{Q*} \\int\\limits_{V} (_{u}N_{i}^{Q}) (_{u}N_{i}^{P}) \\dd{V} u^{P} \\\\\n&- \\delta u^{Q*}\\int\\limits_{V}\\ _{u}N_{i}^Q f_i \\dd{V} - \\delta u^{Q*}\\int\\limits_{S}\\ _{u}N_{i}^Q t_i \\dd{S} \n+ \\delta \\theta^{Q*} \\int\\limits_{V} (_{\\kappa}B_{i}^{Q}) (D_{ij}) (_{\\kappa}B_{j}^{P}) \\dd{V} \\theta^{P} \\\\\n&-\\delta \\theta^{Q*}\\int\\limits_{S}\\ _{\\theta}N_i^Q m_i \\dd{S} + \\delta s^{Q*} \\int\\limits_{V} (_{s}N_{k}^{Q}) (_{\\nabla}B_{k}^{P}) \\dd{V} u^{P}\n+ \\delta u^{Q*} \\int\\limits_{V} (_{\\nabla}B_{k}^{Q}) (_{s}N_{k}^{P}) \\dd{V} s^{P} \\\\\n&-\\delta s^{Q*} \\int\\limits_{V} 2 (_{s}N_{k}^{Q}) (_{\\theta}N_{k}^{P}) \\dd{V} \\theta^{P} -\\delta \\theta^{Q*} \\int\\limits_{V} 2(_{\\theta}N_{k}^{Q}) (_{s}N_{k}^{P}) \\dd{V} s^{P} = 0 \\, .\n\\end{split}\n\\label{eq:discrete_PVW}\n\\end{equation}\n\nThe explicit form of the interpolators defined above is given in the appendix.\n\n\\subsection{Discrete equilibrium equations}\nFrom the arbitrariness in the variations \\(\\delta u^Q\\) , \\(\\delta \\theta^Q\\)\nand \\(\\delta s^Q\\) in \\eqref{eq:discrete_PVW} it follows that:\n\\begin{equation*}\n\\begin{split}\n&\\int\\limits_{V} (_{e}B_{ij}^{Q}) (C_{ijkl}) (_{e}B_{kl}^{P}) \\dd{V} u^P\n - \\int\\limits_{V} (_{u}N_{i}^{Q}) (_{u}N_{i}^{P}) \\dd{V} u^P\n - \\int\\limits_{V}\\ _{u}N_{i}^Q f_i \\dd{V}\n - \\int\\limits_{S}\\ _{u}N_{i}^Q t_i \\dd{S} = 0\\,, \\\\\n&\\int\\limits_{V} (_{\\kappa}B_{i}^{Q}) (D_{ij}) (_{\\kappa}B_{j}^{P}) \\dd{V} \\theta^P \n - \\int\\limits_{S}\\ _{\\theta}N_i^Q m_i \\dd{S}\n - \\int\\limits_{V} 2(_{\\theta}N_{k}^{Q}) (_{s}N_{k}^{P}) \\dd{V} s^P=0\\, , \\\\\n&\\int\\limits_{V} (_{s}N_{k}^{Q}) (_{\\nabla}B_{k}^{P}) \\dd{V} u^P\n - \\int\\limits_{V} 2 (_{s}N_{k}^{Q}) (_{\\theta}N_{k}^{P}) \\dd{V} \\theta^P = 0\\, ,\n\\end{split}\n\\end{equation*}\nwhich can be written in the standard finite element form for dynamic equilibrium\n\\begin{equation}\n\\begin{bmatrix}\nK_{uu}^{QP} &0 &K_{us}^{QP}\\\\\n0 &K_{\\theta\\theta}^{QP} &-K_{\\theta s}^{QP}\\\\\nK_{s u}^{QP} &-K_{s\\theta}^{QP} &0\n\\end{bmatrix}\n\\begin{Bmatrix}\nu^P\\\\\n\\theta^P\\\\\ns^P\n\\end{Bmatrix}\n=\n\\omega^2\\begin{bmatrix}\n M_{uu}^{QP} &0 &0\\\\\n0 &0 &0\\\\\n0 &0 &0\n\\end{bmatrix}\n\\begin{Bmatrix}\nu^P\\\\\n\\theta^P\\\\\ns^P\n\\end{Bmatrix}\n+ \\begin{Bmatrix}\nF_u^Q\\\\\nm_\\theta^Q\\,\\\\\n0\\end{Bmatrix}\n\\label{eq:mat_fem}\n\\end{equation}\nwhere the individual terms are defined as\n\\begin{align*}\nK_{uu}^{QP} &= \\int\\limits_{V} (_{e}B_{ij}^{Q}) (C_{ijkl}) (_{e}B_{kl}^{P}) \\dd{V}\\, ,\n&M_{uu}^{QP} &= \\rho \\omega ^2 \\int\\limits_{V} (_{u}N_{i}^{Q}) (_{u}N_{i}^{p}) \\dd{V}\\, ,\\\\\nK_{u s}^{QP} &= \\int\\limits_{V} (_{\\nabla}B_{k}^{Q}) (_{s}N_{k}^{P}) \\dd{V}\\, ,\n&F_{u }^{Q} &= \\int\\limits_{V}\\ _{u}N_{i}^Q f_i \\dd{V} + \\int\\limits_{S}\\ _{u}N_{i}^Q t_i \\dd{S}\\, ,\\\\\nK_{\\theta \\theta}^{QP} &= \\int\\limits_{V} (_{\\kappa}B_{i}^{Q}) (D_{ij}) (_{\\kappa}B_{j}^{P}) \\dd{V}\\, ,\n&K_{\\theta s}^{QP} &= \\int\\limits_{V} 2(_{\\theta}N_{k}^{Q}) (_{s}N_{k}^{P}) \\dd{V}\\, ,\\\\\nm_{\\theta}^{Q} &= \\int\\limits_{S}\\ _{\\theta}N_i^Q m_i \\dd{S}\\, ,\n&K_{s u}^{QP} &= \\int\\limits_{V} (_{s}N_{k}^{Q}) (_{\\nabla}B_{k}^{P}) \\dd{V}\\, ,\\\\\nK_{s \\theta}^{QP} &= \\int\\limits_{V} 2 (_{s}N_{k}^{Q}) (_{\\theta}N_{k}^{P}) \\dd{V}\\, .\n\\end{align*}\n\n\\Cref{eq:mat_fem} can be rewritten in the following set of equilibrium equations\nin terms of nodal forces and couples\n\\begin{equation}\n\\begin{aligned}\nf_{(\\sigma)}^Q + f_{s}^Q - f_I^Q - T^Q &= 0\\, ,\\\\\nm_{\\mu}^Q + m_{s}^Q - q^Q &= 0\\, ,\\\\\ns(\\theta - \\hat{\\theta}) = 0\\, ,\n\\end{aligned}\n\\label{eq:discrete_balance}\n\\end{equation}\nwhere the subindex $(\\sigma)$ refers to the symmetric part of the stress tensor,\nand \\(I\\) to inertial forces. Notice that we do not have an\ninertial term for the second equation as is the case for the micropolar model\n\\citep{guarin2020}. We also have a third equation reflecting the kinematic\nrestriction, between the rotation \\(\\theta\\) and the introduced degree of\nfreedom \\(\\hat{\\theta}\\), imposed via the Lagrange-multiplier term \\(s\\) in\neach element.\n\nWhen using a Lagrange multiplier formulation as in\n\\eqref{eq:weak_form_with_multipliers} the equations are still self-adjoint, as\ncan be seen in the structure of \\eqref{eq:mat_fem}. Nevertheless, the stiffness\nmatrix is indefinite and the solution of the problem represents a saddle-point\ninstead of a minimum \\citep{arnold_mixed_1990, darrall2014}.\n\n\\subsection{Eigenvalue problem}\nIn finding the dispersion relations, we are interested in the free wave motion\nin the media. This leads to the following eigenvalue problem\n\\begin{equation}\n[K]\\{U\\} = \\omega^2 [M]\\{U\\}\n\\label{eq:eigenvalue_fem}\n\\end{equation}\nwith\n\\[[K] = \\begin{bmatrix}\nK_{uu}^{QP} &0 &K_{us}^{QP}\\\\\n0 &K_{\\theta\\theta}^{QP} &-K_{\\theta s}^{QP}\\\\\nK_{s u}^{QP} &-K_{s\\theta}^{QP} &0\n\\end{bmatrix} ,\\,\n[M] = \\begin{bmatrix}\nM_{uu}^{QP} &0 &0\\\\\n0 &0 &0\\\\\n0 &0 &0\n\\end{bmatrix}, \\,\n\\{U\\} = \n\\begin{Bmatrix}\nu^P\\\\\n\\theta^P\\\\\ns^P\n\\end{Bmatrix}\\, .\\]\n\nIn \\eqref{eq:eigenvalue_fem} Bloch-periodic boundary conditions are yet to be\nimposed. This can be done in two ways \\citep{valencia_uel_2019}: (i) modifying\nthe connectivity of the elements; and (ii) assembling the matrices without\nconsidering boundary conditions and impose the Bloch-periodicity through\nrow\/column operations. In this work, we follow the second approach as it\nrequires the stiffness and mass matrices to be assembled once and the\ntransformation matrices are computed for every wavenumber in the first\nBrillouin zone. This process results in the following eigenvalue problem\n\\begin{equation}\n[K_R(\\vb{k})]\\{U\\} = \\omega^2 [M_R(\\vb{k})]\\{U\\}\n\\label{eq:reduced_eigenvalue_fem}\n\\end{equation}\nwith\n\\[[K_R(\\vb{k})] = [T(\\vb{k})^H K T(\\vb{k})]\\, ,\\quad [M_R(\\vb{k})] = [T(\\vb{k})^H M T(\\vb{k})]\\, ,\\]\nwhere \\([T(\\vb{k})]\\) represents the transformation matrix for a given \\(\\vb{k}\\),\nand the \\([T^H]\\) refers to the Hermitian transpose of \\([T]\\). For an explicit\nform for the matrices \\([T]\\) refer to \\cite{hussein2014dynamics} or\n\\cite{guarin2012_msc}.\n\nWe conducted the implementation on top of the in-house finite element code\nSolidsPy \\citep{solidspy} and used SciPy to solve the eigenvalue problem\n\\citep{scipy}. To take advantage of the sparsity of the matrices the problem\nshould be written as matrix-vector multiplications, such as\n\\begin{align*}\n\\{x\\} = [T]\\{U\\}\\, ,\\\\\n\\{y\\} = [K]\\{x\\}\\, ,\\\\\n\\{z\\} = [T^H] \\{y\\}\\, ,\n\\end{align*}\nwith $\\{z\\}$ representing the image of the linear operator \\([K_R]\\) over\n\\(\\{U\\}\\). The same procedure can be applied for the right-hand side of\n\\eqref{eq:reduced_eigenvalue_fem}.\n\nThe Lagrange-multiplier approach represents a saddle-point instead of a minimization\nproblem \\citep{arnold_mixed_1990}. This can be seen in the structure of the stiffness\nmatrix obtained in equation \\eqref{eq:eigenvalue_fem}. Furthermore the mass matrix\nis not positive definite anymore. This structure for the eigenvalue problem requires\nthe use of a specific solver such as the LOBPCG method \\citep{knyazev2001}\ninstead of the classical Arnoldi method \\citep{book:arpack}.\n\n\\section{Results: Dispersion relations for C-CST cellular materials}\nIn this section we conduct a series of dispersion analyses intended to show the\neffectiveness of our mixed finite element implementation of the C-CST material\nmodel in predicting the correct wave propagation properties of the material. All\nthe dispersion graphs use the dimensionless frequency\n\\begin{equation}\n\\Omega = \\frac{2d\\omega}{c_2}\\, ,\n\\end{equation}\nfor the vertical axis, where \\(2d\\) is the dimension of the unit cell and\n\\(c_2^2 = \\mu\/\\rho\\) is the speed of the shear wave for a classical elastic material.\nThe Poisson ratio for all the simulations is \\(\\nu = 1\/4\\).\n\nAs a first instance we find the response of a homogeneous periodic material\nwhich has also a closed form solution. We will then continue to study a second\nprototypical example corresponding to a homogeneous material with a circular\npore. These two problems exhibit two different levels of dispersive behavior. In\nthe homogeneous material cell, dispersion is due to the kinematic enrichment of\nthe model associated to the length scale parameter, while in the porous material\nmodel additional dispersion arises due to the explicit microstructural feature.\n\n\\subsection{Homogeneous material}\nAs a test of accuracy and effectiveness of our implementation we consider the\ncase of a homogeneous material cell with the same mechanical properties of the\nmaterial reported previously and with closed form dispersion relations from\n\\eqref{eq:dispersion} and \\eqref{eq:dispersion_SV}. In this model microstructural\neffects are introduced through the material length parameter \\(\\ell\\).\nRecall that \\(\\ell^2\\) is defined by the ratio \\(\\frac{\\eta}{\\mu}\\) where\n\\(\\eta\\) is the curvature-couple-stress module while \\(\\mu\\) is the shear\nmodulus from Cauchy elasticity. The results in terms of the resulting band\nstructure are shown in \\cref{fig:homogeneous}, where we used a\n\\(16\\times16\\) mesh and \\(\\ell^2\/d^2 = 3\/8\\). For a conceptual description of\nthe reciprocal space and a guide on how to interpret the results in a Bloch\nanalysis the reader is referred to \\cite{Valencia_periodic}. Note that this set\nof results is directly comparable with the curves from the closed form solutions\nfrom \\cref{fig:dispersion_analytic}. Since the material is isotropic there are\nno directional effects and, as discussed previously, the only difference between\nthis model and the result from classical elasticity is the dispersive behavior\nof the shear wave. In contrast with the micropolar model \\citep{guarin2020},\nthe present C-CST model does not exhibit additional rotational waves.\n\\begin{figure}[H]\n\\centering\n\\includegraphics[height=2.8 in]{img\/comparison_analytical_mod.pdf}\n\\caption{Dispersion relations for a homogeneous material model. Solid lines\nrepresent FEM results while markers correspond to the analytic solution.\nTriangular and filled-dots describe the P and SV wave modes, respectively.}\n\\label{fig:homogeneous}\n\\end{figure}\n\n\\Cref{fig:homog_lengths} shows the results for the same\nmaterial cell but now we have considered 4 different values of the length scale\nparameter corresponding to \\(\\ell\/d \\in [0.01, 0.1, 1, 10]\\). The mesh in this\ncase is \\(16\\times16\\). Notice that, as expected, the increasing value of this\nparameter only affects the dispersive response of the shear waves while the\nP-waves retain their classical non-dispersive behavior. As seen in\n\\eqref{eq:dispersion_SV} the dispersion increases for higher values of \\(\\ell\\)\ndue to the factor \\(\\sqrt{1 + k^2 \\ell^2}\\) in the dispersion relation. This\nbehavior is closely followed by the numerical results presented in\n\\cref{fig:homog_lengths}.\n\\begin{figure}[H]\n\\centering\n\\includegraphics[width=5.5 in]{img\/homogeneous_lengths.pdf}\n\\caption{Dispersion relations for a homogeneous material model with varying\nlength scales, \\(\\ell\/d \\in [0.01, 0.1, 1, 10]\\). Solid lines represent FEM\nresults while markers correspond to analytic solutions. As expected, for\nincreasing \\(\\ell\/d\\) the S-wave presents more dispersion due to the factor\n\\(\\sqrt{1 + k^2 \\ell^2}\\), as presented in \\eqref{eq:dispersion_SV}.}\n\\label{fig:homog_lengths}\n\\end{figure}\n\nAs an additional verification, we also tested the convergence\nin the calculation of the dispersion relations after considering the first 8\nmodes for a sequence of meshes of \\(1\\times1\\), \\(2\\times2\\), \\(4\\times4\\),\nand \\(8\\times8\\) elements for \\(\\ell^2\/d^2 = 3\/8\\). The error in the eigenvalue\ncomputation was measured according to\n\\[e = \\frac{\\Vert\\vb*{\\omega}_\\text{ref} - \n\\vb*{\\omega}_h\\Vert_2}{\\Vert\\vb*{\\omega}_\\text{ref}\\Vert_2}\\, ,\\]\nwhere \\(\\vb*{\\omega}_h\\) is the set of eigenvalues (dispersion relation) for a \nmesh of characteristic element size \\(h\\) and \\(\\vb*{\\omega}_\\text{ref}\\) is the \nsolution corresponding to the finer \\(16\\times16\\) elements mesh, which has been\ntaken as reference. The results for this sequence, together with the variation\nin the error parameter, are displayed in \\cref{fig:convergence}. The estimated \nconvergence rate for the eigenvalues is 2.32. We see that when\nwe refine the mesh it can reproduce the dispersion curves better for higher\nfrequencies. There are still some differences between the \\(8\\times8\\) and\n\\(16\\times16\\) meshes around the dimensionless frequency of 15 but these\ndifferences will disappear with further refinement. Nevertheless, opposed to\nwhat happens in classical continua we would need more points per wavelength\nevery time that we want to increase the maximum frequency. This is due to the\ninherent dispersive behavior of S-waves as can be seen in equation\n\\eqref{eq:dispersion_SV}. Thus, we would expect to need more than 10 points per\nwavelength, customary for finite element methods, or 5, customary for spectral\nelement methods \\citep{komatitsch1999introduction, ainsworth2009dispersive,\nguarin2015}. Again, we should emphasize the dispersive nature of the SV-waves,\nwhile the P-waves remain non-dispersive.\n\n\n\\begin{figure}[H]\n\\centering\n\\includegraphics[width=5.75 in]{img\/convergence.pdf}\n\\caption{Convergence of the first 8 modes in the dispersion relations at\n\\(\\ell^2\/d^2 = 3\/8\\) for a sequence of meshes with: \\(1\\times1\\), \\(2\\times2\\),\n\\(4\\times4\\), and \\(8\\times8\\) elements --- presented as solid blue lines in the \nbackground. The results are compared with a mesh that has \\(16\\times16\\)\nelements --- presented as dots in the foreground. The estimated convergence rate\nfor the eigenvalues in the 2-norm is 2.32.}\n\\label{fig:convergence}\n\\end{figure}\n\n\n\\subsection{Dispersion in cellular material with a circular pore}\nNow, we consider a C-CST composite material cell configured by a circular pore\nembedded in a homogeneous matrix. The presence of the pore provides the model\nwith a second length-scale due to the microstructure, in addition to the one\nintroduced by the material length-scale parameter \\(\\ell\\). For\nillustration, we assume a pore diameter \\(a\\) that is half the cell length\n(i.e., \\(a=d\\)). Thi is equivalent to a porosity of \\(\\pi\/16\\) or approximately\n0.196, which is kept fixed as we modify the size of the unit cell to control\n\\(\\ell\/a\\).\n\nThe resulting dispersion curves for this cellular material\nwith four different length scale ratios are shown in \\cref{fig:pore}\nwith \\(\\ell\/a = [0.01, 0.1, 1, 10]\\). In contrast to the results from the fully\nhomogeneous material cell, the presence of the circular pore introduces\nscattering effects inside each cell and the composite shows much more\ncomplicated elastodynamic behavior. Most importantly, however, the dispersion\ncurves become more regular with increased \\(\\ell\/a\\) and partial bandgaps\nopen up along the \\(\\Gamma M\\) and \\(X\\Gamma\\) directions, especially for\n\\(\\ell\/a = 1\\) and \\(\\ell\/a = 10\\). This type of band structure is not\nobserved for classical elastodynamic cells with a similar geometric periodicity,\nwhich would exhibit behavior close to that obtained here with \\(\\ell\/a = 0.01\\).\nIn fact, as \\(\\ell\/a \\rightarrow 0\\), C-CST theory recovers the classical\nresult.\n\\begin{figure}[H]\n\\centering\n\\includegraphics[width=5.5 in]{img\/cpore_lengths_mod.pdf}\n\\caption{Dispersion for a cellular material with circular pores for varying \nlength scales, $\\ell\/a = [0.01, 0.1, 1, 10]$.}\n\\label{fig:pore}\n\\end{figure}\n\nFurthermore, it is important to note that this interesting band structure is\nobtained with a consistent continuum mechanics formulation, which requires only\na single additional material parameter, \\(\\ell\\), beyond those needed in the\nclassical elastic case. When this intrinsic length scale is on the order of the hole\ndiameter, the dispersive SV wave has a fundamental group velocity for the cell\nwith dimension \\(2d\\) approximately equal to the group (and phase) velocity of\nthe non-dispersive P wave, which allows the SV and P branches to follow a\nsimilar path, causing band gaps to open. This behavior, which is not seen in\nclassical elastodynamics, occurs in the regions near to where the two branches\nintersect in \\cref{fig:dispersion_analytic}. Consequently, under C-CST,\nthese band gaps will originate whenever the size of the cellular structure is\ntuned to the material length scale, a potentially significant phenomenon that\nhas not been recognized previously. With further tuning of the porosity level\nand cell size, it may be possible to achieve even a complete bandgap at\nrelatively low non-dimensional frequency \\(\\Omega\\).\n\n\n\\section{Conclusions}\n\nThe present work incorporates several innovative aspects. First of all, we\nhave developed a novel frequency domain correlated action principle for the\nconsistent couple stress theory (C-CST) of \\cite{hadjesfandiari2011couple} and\nused that to extend the Lagrange multiplier finite element algorithm of\n\\cite{darrall2014} to study periodic cellular materials through Floquet-Bloch\ntheory from solid state physics. Particularly, we have addressed the imposition\nof extended Bloch boundary conditions for this material model where in\naddition to force tractions and displacements there are also couple tractions\nand rotations. Secondly, we also discussed numerical aspects related to the\nsolution of the wavenumber dependent generalized eigenvalue problem resulting\nfrom the imposition of the Bloch periodic boundary conditions, overcoming\ncomplications arising from the inclusion of Lagrange multipliers and a\nnon-positive definite mass matrix. The implementation was\nshown to give accurate results for homogeneous and porous unit cells and for\nvarying couple stress material length-scale parameters.\n\nThe analysis of the first cell was used mainly to test the correctness of our\nimplementation as this material has a closed-form dispersion relation. The\nalgorithm was shown to correctly capture the non-dispersive P-wave as well as\nthe dispersive SV-wave. This analysis was complemented by a convergence analysis\nwith four different meshes of increasing refinement for the material cell. The\nobserved convergence rate shows that the Lagrange multiplier algorithm is\neffective in maintaining continuity by imposing the newly introduced kinematic\nconstraint implicit in the mean curvature tensor definition. As the final\ncontribution, we have discovered the interesting bandgap structure of a material\ncell with a circular pore embedded in a homogeneous matrix, which reveals the\nappearance of bandgaps introduced by the kinematic features of C-CST and the\ndispersive behavior of the SV-waves defined in terms of the microstructural\nlength scale parameter.\n\nFrom a general perspective, C-CST is a true size-dependent continuum theory,\nwhich is intended here for periodic elastic material cells at scales for which a\ncontinuum representation is appropriate. From the results shown in the paper,\nC-CST becomes important when the size of the cell in on the order of the\nintrinsic length scale parameter or smaller. For larger cells, the classical\ntheory can be used instead. On the other hand, micropolar theory disconnects the\nrotational field from the displacements, which can lead to approximations that\nmay or may not be physical.\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\tConvolutional Neural Networks (CNNs) have been successfully used in various tasks with Euclidean structure data in recent years. For non-Euclidean structure datasets, graph convolutional networks (GCNs) use the same idea to extract the topological structure information.\n\t\n\t\\cite{bruna2014spectral} firstly proposed two GCN models with spectral and spatial construction respectively. The spectral model using a Laplacian matrix to aggregate neighborhood information of each node in a graph. The spatial model partition graph into clustering and update them by aggregating function. In order to extract deeper features, model based on CNNs usually deepen the model's layers. While in GCNs, Deepening layers will cause a lot of problems. In spectral construction method, too many layers lead to \\textbf{``over smoothing''}\\cite{li2018deeper,zhou2018graph}: the features of nodes in graph will tend to be the same. In spatial construction method, it will cause exponential growth of the number of sampled neighbor nodes, called \\textbf{``neighbor explosion''}. Node sampling and layer sampling\\cite{hamilton2018inductive,chen2018fastgcn,gao2018large,huang2018adaptive,ying2018hierarchical} were proposed to handle this problem, but due to incomplete sampling, the inaccuracy of nodes' representation accumulates errors between layers.\n\t\n\tTo address these two problems, in this paper, we propose a model called PathSAGE, which can learn high-order topological information by expanding the receptive field.\n\tFirstly, We design a path sampling technique based on random walk to sample paths starting from central nodes with different lengths, then the sequences of paths are fed into Transformer encoder\\cite{vaswani2017attention}, which can extract the semantic and distance information in sequence effectively. As shown in Fig. 1, We view the sequences of paths from the tail nodes to the central node as the central node' neighbors. Secondly, following this idea, we take the average of paths of the same length as the representation of the central node in this level of reception field. Finally, after concatenating the aggregated features of paths with different lengths, the final representation of the central node are used for downstream tasks.\n\t\n\tFor the two problems mentioned earlier, on the one hand, the aggregation of the central node only perform once in training of a sample, which never cased ``over-smoothing''. On the other hand, all the paths were sampled with a fixed number for representing the central node in our model, instead of recursively sampling exponentially growing neighbor nodes. And each path only contributes to the central node, we do not need to calculate and store the temporary representations of nodes from the middle layer. Furthermore, it prevents the error propagation caused by incomplete sampling. \n\t\n\tOur contribution can be summarized in three points:\n\t\n\t\\begin{itemize}\n\t\t\n\t\t\\item We utilize the path sampling to take place of the node sampling to avoid error accumulation caused by incomplete node sampling.\n\t\t\n\t\t\\item We propose and evaluate our model Path-SAGE to solve the existing ``neighour explosion'' and ``over-smoothing'' problems. the model can capture richer and more diverse patterns around the central node with only one layer of structure to the central node.\n\t\t\n\t\t\\item We evaluate our model on three inductive learning tasks, and it reaches the state-of-the-art performance on two of them. We analyze the attention weights of Transformer encoder and detect some patterns in the attention mechanism, which can further illustrate how the model works.\n\t\t\n\t\\end{itemize}\n\t\\section{Related work}\n\tGNNs model was initially proposed by \\cite{bruna2014spectral}, and the convolution operation in the traditional Euclidean structure was introduced into the graph network with the non-Euclidean structure in this article. They \\cite{bruna2014spectral} divided GNN into two construction methods: spectral construction and spatial construction. Subsequently, many studies are carried out around these two aspects. \n\t\n\tIn spectral construction, \\cite{defferrard2017convolutional} used Chebyshev polynomials with learnable parameters to approximate a smooth filter in the spectral domain, which improves computation efficiency significantly. \\cite{kipf2017semisupervised} further reduced the computational cost through local first-order approximation. In spatial construction, MoNet was proposed in \\cite{monti2017geometric}, developing the GCN model by defining a different weighted sum of the convolution operation, using the weighted sum of the nodes as the central node feature instead of the average value. \\cite{hamilton2018inductive} attempted various aggregator to gather the features of neighbor nodes. \\cite{velivckovic2017graph,zhang2018gaan} defined the convolution operation with a self-attention mechanism between the central node and neighbor nodes, \\cite{liu2019geniepath} brought the LSTM(Long Short Term Memory networks)\\cite{hochreiter1997long} from NLP (Natural Language Processing) to GNN, and built an adaptive depth structure by applying the memory gates of LSTM. \\cite{yang2021spagan} sampled the shortest path based on their attention score to the central node and combined them in a mixed way. PPNP and APPNP\\cite{klicpera2018predict} used the relationship between GCN and PageRank to derive an improved communication scheme based on personalized PageRank. \n\t\n\tWith the scale of graph data increasing, the full-batch algorithm is no longer applicable, and the mini-batch algorithm using stochastic gradient descent is applied. In recent years, some studies based on different sampling strategies have been proposed. \\cite{hamilton2018inductive} tried random neighbor node sampling for the first time to limit large amounts of nodes caused by recursive node sampling. Layer-wise sampling techniques was applied in \\cite{chen2018fastgcn,gao2018large,huang2018adaptive,ying2018hierarchical}, which only consider fixed number of neighbors in the graph to avoid ``neighbor explosion''. \\cite{chen2018fastgcn,chen2018stochastic} used a control variate-based algorithm that can achieve good convergence by reducing the approximate variance. Besides, subgraph sampling was first introduced by \\cite{zeng2019accurate}, which used a probability distribution based on degree. \\cite{chiang2019cluster} performed clustering decomposition on the graph before the training phase and randomly sampled a subgraph as a training batch at every step. \\cite{zeng2020graphsaint} sampled subgraph during training based on node and edge importance. All these sampling techniques are based on GCNs, how to extract deeper topology information from large graph datasets is still a problem.\n\t\\section{Proposed method}\n\tIn this section, we present the PathSAGE. Firstly, the sampling algorithm is introduced in 3.1. Secondly, we detail the aggregator in 3.2. Finally, we discuss the difference between our model and the related models.\n\t\\begin{algorithm} \n\t\t\\SetAlgoLined\n\t\t\\DontPrintSemicolon\n\t\t\\caption{Random Path Sampling} \n\t\t\t\\KwIn{garph $G(V,E)$\\\n\t\t\tcentral node $c$\\\\\n\t\t\tsample depth $s$\\\\\n\t\t\tsample num each length $L =\\{n_1, n_2, \u2026, n_s\\}$\\} }\n\t\t\t\\KwOut{path sequcences with different length\\\\ \n\t\t\t$\\{\\bm{P}_1, \\bm{P}_2, ..., \\bm{P}_s\\}$}\n\t\t\t\\SetKwFunction{Fun}{Ramdom Path Sampling}\n\t\t\t\\Fun{$G$, $c$, $s$, $L$}{\n\t\t\t\t\n\t\t\t\\ForEach{$l = 0 \\to s$}{\n\t\t\t$i\\gets0$;\n\t\t\t\n\t\t\t\\While{$i < n_l$}{\n\t\t\t $u\\gets c$;\n\t\t\t \n\t\t\t $P\\gets \\{u\\}$;\n\t\t\t \n\t\t\t\\For{$j = 0 \\to l$}{\n\t\t\t $u \\gets$ Node randomly selected from $u$'s neighbors;\n\t\t\t \n\t\t\t $P\\gets P \\cup \\{u\\}$}\n\t\t \n\t\t\t $\\bm{P}_l \\gets \\bm{P}_l \\cup \\{P\\}$;\n\t\t\t \n\t\t\t $i \\gets i+1$;}\n\t\t \n\t\t\t}\n\t\t\\Return $\\{\\bm{P}_1, \\bm{P}_2, ..., \\bm{P}_s\\}$\n\t\t}\n\t\\end{algorithm}\n\t\\subsection{Path Sampling}\n\tExcept deepening model, another way to expand receptive field of CNNs is to increase the size of convolution kernel. Following this idea, we sample the node sequences starting from central node and regard these paths as the context of the central node. Therefore, the receptive field can be expanded by extending these paths.\n\t\n\tWe use a straightforward random sampling algorithm based on random walk, shown in \\textbf{Algorithm 1}: for a central node, a sampling starts from it and randomly selects a neighbor node each time until reaching the preset path length, and multiple path sequences for the corresponding central node of various lengths can be obtained in this way, constituting a training sample.\n\n\t\n\t\\begin{figure*}[htb]\n\t\t\\begin{center}\n\t\t\t\\includegraphics[height=4.5cm]{.\/fig1.pdf}\n\t\t\\end{center}\n\t\t\\caption{Mechanism of sampling and aggregation. (a) and (b) are two different possible training samples with the same central node. Paths with same colors are with same lengths same length and share the same aggregators. }\n\t\n\t\t\\label{fig1}\n\t\\end{figure*}\n\t\n\t\\begin{figure}[htbp]\n\t\t\\begin{center}\n\t\t\t\\includegraphics[height=4.5cm]{.\/fig2.pdf}\n\t\t\\end{center}\n\t\t\\caption{Structure of first aggregation in aggregator. A specific sampled path's features adds the position embedding, through Transformer encoder layers to fuse the information.\n\t\t}\n\t\t\\label{fig2}\n\t\\end{figure}\n\t\n\t\\subsection{Aggregator}\n\tThere are two aggregations in the aggregator: the first one aggregate path sequences; the second one aggregate different paths as the final representation of the central node. \t\n\tFor the first aggregation, We formulate each path sequence as:\n\t\\begin{equation}\n\t\t{\\bm{P}}_{ij}^{}\\; = \\;\\{ {\\vec a_{ij}^1,\\;\\vec a_{ij}^2,\\;\\vec a_{ij}^3,\\;...,\\;\\vec a_{ij}^j\\;} \\}\n\t\\end{equation}\n\n\t\n\tTo utilize the position information of path sequence, we define positions of nodes in paths as their distances to the central nodes. As same as Transformer, we add nodes' features and their positional vectors together as the input of the structure:\n\t\\begin{equation}\\label{eq6}\n\t\t\\bm{\\tilde{P}}_{ij}\\; = \\;\\{ {a_{ij}^1\\; + \\;pos\\_emb\\left( 1 \\right),...,\\;a_{ij}^j\\; + \\;pos\\_emb\\left( j \\right)} \\}\n\t\\end{equation}\n\tFor the $pos\\_emb(\\cdot)$, we use $sine$ and $cosine$ functions of different frequencies to generate the positional vectors. In each dimension, the positional embedding is:\n\t\n\t\\begin{equation}\\label{eq7}\n\t\tpos\\_emb(p)_{2i} = sin(\\frac{{p}}{{{{10000}^{\\frac{{2i}}{{{d}}}}}}})\n\t\\end{equation}\n\t\n\t\\begin{equation}\\label{eq8}\n\t\tpos\\_emb(p)_{2i + 1} = cos(\\frac{{p}}{{{{10000}^{\\frac{{2i}}{{d}}}}}})\n\t\\end{equation}\n\tWhere $p$ is the position in the sequences, $2i$ and $2i+1$ mean the odd and even dimensions of position embedding, and $d$ is the number of features. \n\t\n\tAfter that, we apply the Transformer encoder on each path: \n\t\n\t\\begin{equation}\\label{eq9}\n\t\t\\bm{\\tilde{P}}_{ij}^{k} = transformer\\_bloc{k^k}(\\bm{\\tilde{P}}_{ij}^{k-1})\n\t\\end{equation}\n\tWhere $k$ means $k$-th Transformer encoder layer. Noted that, in a $m$-layer Transformer encoder, the output of the last layer $\\bm{\\tilde{P}}_{ij}^{m}$ is a sequence of features. We only take the output at position 0 as the final representation of the path sequence.\n\t\n\t\\begin{equation}\\label{eq10}\n\t\t\\tilde{P}_{ij}' = \\left[\\bm{\\tilde{P}}_{ij}^{m}\\right]_{0}\n\t\\end{equation}\t\n\t\n\tFollowing these equations (\\ref{eq6})-(\\ref{eq10}), we can obtain representations of paths with different length to the central node, as shown in Fig. 1. In the second aggregation, we apply an average pooling layer to aggregate the central node paths. Then we concatenate all the path representation and apply a feed-forward layer with non-linearity to fuse these features. The final output of a central node $C'$ is computed as following:\n\t\n\t\\begin{equation}\n\t\t\\begin{array}{c}\n\t\t\tC = \\;concat\\left( {{C_1},\\;...,\\;{C_s}} \\right)\\\\\n\t\t\t\\\\\n\t\t\twhere\\;{C_i}\\; = \\;Average\\left( \\tilde{P}_{i1}',\\;...,\\;\\tilde{P}_{in}' \\right)\n\t\t\\end{array}\n\t\\end{equation}\n\t\n\t\\begin{equation}\n\t\t{C'}\\; = \\;max\\left( {0,\\;C{\\bm{W}_1}\\; + \\,{b_1}} \\right){\\bm{W}_2}\\; + \\;{b_2}\n\t\\end{equation}\n\twhere $s$ denotes the sample depth, $n$ denote the number of paths sampled in\n\t\n\t\\subsection{Comparisons to related work}\n\n\t\\begin{itemize}\n\t\t\\item We introduce a sequence transduction model into our structure, but it is distinct from related work based on these models. LSTM\\cite{hochreiter1997long} was also used in GraphSAGE\\cite{zeng2020graphsaint} to aggregate node's features, it is very sensitive to the order of the input sequence. In contrast, neighbor nodes are disordered, and the authors have rectified it by consistently feeding randomly-ordered sequences to the LSTM. In our model, paths to the central node already have orders, which are naturally suitable for sequence processing model. \n\t\t\\item The difference between our attention mechanism and GAT is that we collect all the mutual attention information of each node in the path. The output of our model is integrated information of entire sequence, instead of only attend to the central node.\n\t\t\\item Path sampling method is also used in SPAGAN\\cite{yang2021spagan}. But each sampled path in SPAGAN is the shortest one to the central node, this sample technique will leads to a high computational overhead and limits the ability of the model to be applied to large graph datasets. By contrast, we sample paths randomly, which greatly save the computation overhead. Besides, in our sample algorithm, the same node may have a different path to the central node, which may help the model acquire more diverse patterns around the central node while saving the computational overhead. \n\t\\end{itemize}\n\t\\section{Experiments}\n\tIn this section, we introduce the datasets and the experiment setting in 4.1 and 4.2 respectively. we present the results of evaluation in 4.3. \n\t\\begin{table}[htbp]\n\t\t\\caption{Summary of inductive learning tasks' statistics.}\n\t\t\\begin{center}\n\t\t\n\t\t\t\\centering\n\t\t\t\\renewcommand{\\arraystretch}{2.2}\n\t\t\n\t\t\t\\setlength{\\tabcolsep}{4mm}{\n\t\t\t\t\\begin{tabular}{|c|c|c|c|}\n\t\t\t\t\t\\hline\n\t\t\t\t\t& \\textbf{Reddit} & \\textbf{Yelp} & \\textbf{Flickr} \\\\\n\t\t\t\t\t\\hline\n\t\t\t\t\ttype & single-label & multi-label & multi-label \\\\\n\t\t\t\t\t\n\t\t\t\t\t\\#Node & 232,965 & 716,847 & 89,250 \\\\\n\t\t\t\t\t\n\t\t\t\t\t\\#Edges & 11,606,919 & 6,977,410 & 899,756 \\\\\n\t\t\t\t\t\n\t\t\t\t\t\\#Features & 602 & 300 & 500 \\\\\n\t\t\t\t\t\n\t\t\t\t\t\\#Classes & 41 & 100 & 7 \\\\\n\t\t\t\t\t\n\t\t\t\t\t\\makecell[c]{Train \/ Val \\\\\/ Test} & \\makecell[c]{66\\% \/ 10\\% \\\\\/ 24\\%} & \t\\makecell[c]{75\\% \/ 10\\% \\\\\/ 15\\%} & \\makecell[c]{50\\% \/ 25\\% \\\\\/ 25\\%} \\\\\n\t\t\t\t\t\\hline\n\t\t\t\\end{tabular} }\n\t\t\\end{center}\n\t\t\\label{tab2}\n\t\\end{table}\n\t\\subsection{Dataset}\n\tThree large-scale datasets are used to evaluate our model: 1) \\textbf{Reddit}\\cite{hamilton2018inductive}: a collection of monthly user interaction networks from the year 2014 for 2046 sub-reddit communities from Reddit that connect users when one has replied to the other, 2) \\textbf{Flickr}: a social network built by forming links between images sharing common metadata from Flickr. 3) \\textbf{Yelp}\\cite{zeng2020graphsaint}: a network that links up businesses with the interaction of its customers. Reddit is a multiclass node-wise classification task; Flickr and Yelp are multilabel classification tasks. \n\tThe detail statistics of these datasets are shown in Tab 1.\n\t\n\t\\begin{table}[t]\n\t\t\n\t\t\\caption{Perforamnce on inductive learning tasks(Micro-F1).}\n\t\t\\begin{center}\n\t\t\n\t\t\t\\centering\n\t\t\t\n\t\t\t\\renewcommand{\\arraystretch}{2}\n\t\t\t\\setlength{\\tabcolsep}{4mm}{\n\t\t\t\t\\begin{tabular}{|c|c|c|c|}\n\t\t\t\t\t\\hline\n\t\t\t\t\t\\textbf{Model} & \\textbf{Reddit} & \\textbf{Yelp} & \\textbf{Flickr} \\\\\n\t\t\t\t\t\\hline\n\t\t\t\t\tGCN & 0.933$\\pm$0.000 & 0.378$\\pm$0.001 & 0.492$\\pm$0.003 \\\\\n\t\t\t\t\t\n\t\t\t\t\tGraphSAGE & 0.953$\\pm$0.001 & 0.634$\\pm$0.006 & 0.501$\\pm$0.013 \\\\\n\t\t\t\t\t\n\t\t\t\t\tFastGCN & 0.924$\\pm$0.001 & 0.265$\\pm$0.053 & 0.504$\\pm$0.001 \\\\\n\t\t\t\t\t\n\t\t\t\t\tS-GCN & 0.964$\\pm$0.001 & 0.640$\\pm$0.002 & 0.482$\\pm$0.003 \\\\\n\t\t\t\t\t\n\t\t\t\t\tAS-GCN & 0.958$\\pm$0.001 & - & 0.504$\\pm$0.002 \\\\\n\t\t\t\t\t\n\t\t\t\t\tClusterGCN & 0.954$\\pm$0.001 & 0.609$\\pm$0.005 & 0.481$\\pm$0.005 \\\\\n\t\t\t\t\t\n\t\t\t\t\tGraphSAINT & 0.966$\\pm$0.001 & \\bf{0.653$\\pm$0.003} & 0.511$\\pm$0.001 \\\\\n\t\t\t\t\t\n\t\t\t\t\tours & \\bf{0.969$\\pm$0.002} & 0.642$\\pm$0.005 & \\bf{0.511$\\pm$0.003} \\\\\n\t\t\t\t\t\\hline\n\t\t\t\\end{tabular} }\n\t\t\t\\label{tab4}\n\t\t\\end{center}\n\t\\end{table}\n\t\\subsection{Experiment setup}\n\t\n\tWe build our model on Pytorch framework\\cite{paszke2017automatic} and construct the Transformer encoder based on UER\\cite{zhao2019uer}. For all tasks, We train the model with Adam SGD optimizer\\cite{kingma2014adam} with learning rate 1e-3 and a learning rate scheduler with 0.1 warmup ratio. We use two layers Transformer, each of which has 8 attention heads to gather the features of paths. The batch size is 32. The dropout layers are applied between each sub-layer in Transformer layer. We use different dropout rates in the output layer and Transformer encoder, which are 0.3 and 0.1 respectively. We set the sampling length of the path ranging from 1 to 8 (depth $s = 8$), and the number of paths sampled in each length are [5, 5, 5, 5, 5, 10, 10, 10]. For multi-label classification task Flickr and Yelp, the final output is obtained through a $sigmoid$ activation, and in Reddit, the final output is obtained through a $softmax$ activation. The hidden dimension is 128 for Reddit and Flickr, 512 for Yelp. \n\t\n\t\\subsection{Result}\n\tWe compare our model with seven state-of-the-art model: GCN\\cite{kipf2017semisupervised}, GraphSAGE\\cite{hamilton2018inductive}, FastGCN\\cite{chen2018fastgcn}, S-GCN\\cite{chen2018stochastic}, AS-GCN\\cite{huang2018adaptive}, ClusterGCN\\cite{chiang2019cluster}, GraphSAINT\\cite{zeng2020graphsaint}. GraphSAGE uses a random node sampling and LSTM aggregator. FastGCN and S-GCN use a control variate-based algorithm that can achieve good convergence by reducing the sampling variance. ClusterGCN performs clustering decomposition on the graph before the training phase and randomly sampling a subgraph as a training-batch at every step. GraphSAINT samples subgraph while training with several sampling methods based on node and edge importance. The evaluation on these tasks uses the $Micro-F1$ metric, and report the mean and confidence interval of the metrics by five runs.\n\t\n\tThe results of inductive learning experiments are shown in Tab 2. As we can see from the table, For Reddit, our model outperforms all the baseline models. For Flickr, we achieve a comparable F1-score with the top-performing model. For Yelp, we surpass most GCN models, second only to GraphSAINT. One hypothesis to explain the difference of the results is: Reddit and Flickr have more training samples and number of nodes features, which makes the attention mechanism have enough data to capture the relationships between nodes.\n\t\n\t\\begin{figure*}[t]\n\t\t\\begin{center}\n\t\t\t\\includegraphics[height=4.7cm]{.\/emp.pdf}\n\t\t\\end{center}\n\t\t\\caption{Detection of attention weights in Transformer that node have the same label in the path sequence may receive higher weights at some attention heads. The heatmap in the upper-left corner is the weight of position 4 of different layers attend to position 2 and attention heads, the right one is the complete heatmap of the selected head. The lower-left picture shows the labels of the nodes in the sequence. }\n\t\t\\label{fig3}\n\t\\end{figure*}\n\t\n\t\\section{Attention Analysis}\n\tWe observed the attention weight of the trained model in the PubMed test set. At some attention heads, we find that nodes with the same labels get very high attention scores to each other. We visualize an example in Fig 3. Central node 0 and node 4 have the same label and receive extremely high attention scores on second attention heads of the first layer. By observing the attention scores for the whole sequence, we can see that this score also occupies a significant share in the sequence. \n\t\n\tThis observation proves that the attention mechanism can successfully capture the information that is helpful to the downstream tasks in the aggregation of paths.\n\t\n\t\\section{Conclusion}\n\tIn this paper, we propose a model---PathSAGE, which can expand the size of receptive field without stacking layers which solved the problem of ``neighbor explosion'' and ``over smoothing''. Regarding all paths as the neighbor nodes of the central node, PathSAGE samples paths starting from the central node based on random walk and aggregates these features by a strong encoder---Transformer. Consequently, the model can obtain more features of neighbor nodes and more patterns around the central node. In our experiment, PathSAGE achieves the state-of-the-art on inductive learning tasks. Our model provides a novel idea to handle large-scale graph data.\n\t\n\t\\section*{Acknowledgments}.\n\tThis work is supported by National Key R\\&D Program of China (No.2017YFB1402405-5), and the Fundamental Research Funds for the Central Universities (No.2020CDCGJSJ0042). The authors thank all anonymous reviewers for their constructive comments.\n\t","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\n\\section{Conclusion}\nIn this paper, we propose a Spatio-Temporal feature auto-selective (STAS) approach that can automatically extract optimal Spatio-Temporal (ST) features hidden in EC for Bias Correcting on Precipitation (BCoP). Experiments on EC benchmark datasets in Eastern China indicate that STAS achieves the highest threat score (TS) on BCop than other 8 algorithms, and has a strong correcting ability in dealing with different degree of precipitation, especially for heavy precipitation. In the future, we will study how to employ ST mechanisms on more complex precipitation scenarios such as squall line, severe convection and thunderstorm. \n\n\n\n\\section{Experiments}\nWe conduct all experiments on time-series EC benchmarks collected from 1) the high-resolution version of the public European Centre (EC) dataset and 2) Meteorological Information Comprehensive Analysis and Process System (MICAPS)~\\cite{luo2006introduction} that can provide the labels of 5 Meteorological Elements (MEs) including precipitation, temperature, pressure, wind, and dew. Our experiments on Bias Correcting on Precipitation (BCoP) mainly contain two parts. The first part compares the Spatio-Temporal feature Auto-Selective (STAS) model with 8 published machine learning (ML) methods. The second part is a set of ablation experiments on STAS. \n\n\\begin{table}[htp]\n\\LARGE\n\\caption{The details of STAS. $N \\bigtimes C \\bigtimes \\ell \\bigtimes(h \\bigtimes w)_{29\\sim3}$ are multi-scale (from 29*29 to 3*3) dimensions of inputs. $(*)_{[.]}$ represent a operator layer and its parameters setting such as filter size from CNN and output scale from Adaptive Pooling (ADP) or Up-Sampling (UpSp). Specifically, the last parameter of Deformable CNN (D-CNN) is spatial dilation rate. Besides, $\\oplus Noisy$ is addition operation of Gaussian noise, and $(\\cdot)$ is shape of outputs in current module along with pipeline $\\rightarrow$.}\n\\smallskip\n\\centering\n\\resizebox{230pt}{25mm}{\n\\smallskip\\begin{tabular}{l|l|l}\n\\hline\n\\multicolumn{3}{c}{{\\bf Inputs:} $N \\bigtimes C \\bigtimes \\ell \\bigtimes(h \\bigtimes w)_{29\\sim3}$} \\\\\n\\hline \n\\multirow{2}*{SFM} & $CNN_{[1\\times1,0]\\rightarrow[3\\times3,0]}$ & $(N\\bigtimes C\\bigtimes (\\ell-1)\\bigtimes (h\\bigtimes w))$ \\\\\n~ & $D-CNN_{[3\\times3,1,0.8]\\rightarrow[3\\times3,1,0.6]} \\rightarrow ADPooling_{[1\\times1]} \\rightarrow FC$ & $\\rightarrow (N\\bigtimes C \\bigtimes(\\ell-1)\\bigtimes(1\\bigtimes1)) \\rightarrow(N)\\bigtimes(\\ell-1) \\Rightarrow \\vct{s}^*$\\\\\n\\hline\n\\hline\n\\multirow{2}*{Encoder} & $CNN_{[1\\times1,0]\\rightarrow[3\\times3,0]} \\rightarrow ADP_{[16\\times16]}$ & $ \\rightarrow (N\\bigtimes C \\bigtimes \\ell \\bigtimes16\\bigtimes16)$ \\\\\n~ & $\\oplus Noisy$ & $\\oplus \\rightarrow (N\\bigtimes C \\bigtimes \\ell \\bigtimes18\\bigtimes18)$ \\\\\n\\hline\n\\hline\n\\multirow{2}*{TFM} & $3DCNN_{[3\\times3,1]} \\rightarrow (ADP_{[1\\times1]}+FC) \\bigtimes 3$ & $\\rightarrow(N\\bigtimes C\\bigtimes18\\bigtimes18) \\rightarrow(N\\bigtimes6) \\bigtimes 3$ \\\\\n~ & $Concat\\rightarrow FC\\bigtimes2$ & $\\rightarrow (N\\bigtimes18) \\rightarrow (N) \\Rightarrow l^{*}$\\\\\n\\hline\n\\hline\nDecoder & $UpSp_{[18\\times18]}$ & $\\rightarrow (N \\bigtimes C\\bigtimes \\ell \\bigtimes 18 \\bigtimes 18)$ \\\\\n\\hline\n\\hline\nRC & $ CNN_{[3\\times3,1]}\\rightarrow ADP_{[1\\times1]} \\rightarrow FC$ & $(N\\bigtimes C\\bigtimes29\\bigtimes29) \\rightarrow (N\\bigtimes C\\bigtimes 1\\bigtimes1) \\rightarrow (N)$ \\\\\n\\hline\n\\multicolumn{3}{c}{{\\bf Output:} $N$} \\\\\n\\hline\n\\end{tabular}\n}\n\\label{details}\n\\end{table} \n\n\\subsection{Datasets and Training Details}\n\\noindent {\\bf EC benchmarks (ECb)} are sliced from a high-resolution version of the public EC dataset~\\cite{berrisford2009era}, only covering Eastern China between ranging from June 1st to August 31st for three years (2016-2018). Concretely, ECb consists of 57 weather features (channels) worldwide selected from 601 meteorological factors by Pearson correlation analysis~\\cite{benesty2009pearson}. Every feature stems from a grid where each pixel in the grid means a specific location. The spatial and temporal resolutions of ECb are about $111km$ per pixel and 6h per time-level respectively. Then ECb is divided into four datasets for different experiments. The first three datasets are 1) ECb only including Tiny rainfall (ECbT), 2) ECb only including Moderate rainfall (ECbM) and 3) ECb only including Heavy rainfall (ECbH) separated into half-open range of precipitation intensity interval ($[0,1mm)$, $[1-10mm)$, $[10mm,+\\infty)$). The last one is ECb Mixed 3 rainfall above (ECbMi) and its mixture ratio of samples is $T:M:H = 9:3:1$.\n\n\\noindent {\\bf Labels} are observations of 5 Meteorological Elements (MEs) from MICAPS in the specific locations in Eastern China every $6h$.\n\n\\noindent {\\bf Training details} are shown in Tab.~\\ref{details}. The table includes the main structure of layers and initial settings of modules in STAS. Besides, the batch size of training is 256 and the testing is 64. We employ Adam as optimizers for all modules in STAS and the learning rate is $1e-4$. Besides, the weight ratios on 5MEs for spatio-temporal losses are set as $2:1:1:1:1$ in which rainfall is 2 and other elements are 1. The uniform scale after upsampling and downsampling is set as $16*16$. For ordinal regression and C3AE, the ranking intervals for precipitation value are set as 0.5 and 1.5 respectively. Specifically, we set a constant standard $1e-3$ as white Gaussian noise. We test our model in every 6 epochs on training and the max epoch of training is 80. Finally, all experiments are conducted in 8 NVIDIA GPUs.\n\\begin{table}[htp]\n\\caption{The 5 criteria between 8 machine learning methods and STAS on ECbMi as in Tab.~\\ref{contrastiveI}. $TS_{0.1}$ is TS score in $PI>0.1$. $TS_{1} | TS_{10}$ is $PI>1$ and $PI>10$ respectively. The ECb forecasts are results of predictions from EC benchmarks themselves. SVR is support vector regression, LR is linear regression, MLP is multilayer perceptron, FCN is full convolutional network, FPN is feature pyramid network, LSTM long short-term memory, OBA is ordinal boosting auto-encoder.}\n\\smallskip\n\\centering\n\\resizebox{240pt}{20mm}{\n\\smallskip\\begin{tabular}{c|c|c|c|c|c}\n\\hline \n &\\multicolumn{5}{|c}{Criteria} \\\\\n\\hline \nMethods& MAE & MAPE & $TS_{0.1}$ & $TS_{1}$ & $TS_{10}$ \\\\\n\\hline \nECb forecasts~\\cite{ran2018evaluation} & 1.76 & 17.09 & 0.41 & 0.31 & 0.19 \\\\\n\\hline\nSVR~\\cite{srivastava2015wrf} & 1.67 & 15.81 & 0.48 & 0.37 & 0.1 \\\\\n\\hline\nLR~\\cite{hamill2012verification} & 1.73 & 16.90 & 0.35 & 0.35 & 0.21 \\\\\n\\hline\nMLP~\\cite{yuan2007calibration} & 1.59 & 15.13 & 0.46 & 0.39 & 0.21 \\\\\n\\hline\nFCN~\\cite{xu2019towards} & 1.26 & 12.30 & 0.49 & 0.48 & 0.24 \\\\\n\\hline\nFPN~\\cite{lin2017feature} & 1.15 & 7.38 & 0.56 & 0.51 & 0.27 \\\\\n\\hline\nLSTM~\\cite{xingjian2015convolutional} & 1.21 & 9.8 & 0.52 & 0.48 & 0.24 \\\\\n\\hline\nOBA~\\cite{xu2019towards} & 1.01 & 8.96 & 0.58 & 0.53 & 0.25 \\\\\n\\hline\nSTAS(ours) & {\\bf 0.98} & {\\bf 5.84} & {\\bf 0.75} & {\\bf 0.69} & {\\bf 0.38} \\\\\n\\hline\n\\end{tabular}\n}\n\\label{contrastiveI}\n\\end{table} \n\\begin{table*}[t]\n\\centering\n\\caption{Ablation experiments conducted on ECbMi using $TS_{1}$ and $MAPE$. $\\surd$ is defined as an existing component in current ablated STAS for every line of Table. Instead, $\\times$ is no this component in the ablation. SFM-MSMs and TFM-MTMs are spatial and temporal meteorological elements modules severally. Besides, D-CNN is deformable CNN and R, T, P, W, and D represents the module of Rainfall, Temperature, Pressure, Wind and Dew respectively.}\n\\fontsize{8}{9}\\selectfont\n \\begin{tabularx}{15.5cm}{@{\\extracolsep{\\fill}}|c|c|ccccc|ccccc|c|c||c|c|}\n \\hline\n \\multirow{2}{*}{SFM} & \\multirow{2}{*}{TFM} & \\multicolumn{5}{c|}{SFM-MSMs} & \\multicolumn{5}{c|}{TFM-MTMs} & \\multirow{2}{*}{D-CNN} & \\multirow{2}{*}{C3AE} & \\multirow{2}{*}{$TS_1$} & \\multirow{2}{*}{$MAPE$} \\\\\n\n\n \n & & R & T & P & W & D & R & T & P & W & D & & & & \\\\\n \\hline\n $\\surd$ & $\\times$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & \\multicolumn{5}{c|}{N\/A} &$\\surd$ & N\/A & 0.60 & 8.58 \\\\\n \\hline\n $\\times$ & $\\surd$ & \\multicolumn{5}{c|}{N\/A} & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ &N\/A & $\\surd$ & 0.64 & 8.01 \\\\\n \\hline\n $\\surd$ & $\\surd$ & $\\surd$ & $\\times$ & $\\times$ & $\\times$ & $\\times$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & 0.62 & 7.97 \\\\\n \\hline \n $\\surd$ & $\\surd$ & $\\times$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & 0.65 & 8.15 \\\\\n \\hline \n $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\times$ & $\\times$ & $\\times$ & $\\times$ & $\\surd$ & $\\surd$ & 0.65 & 7.03 \\\\ \n \\hline \n $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\times$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & 0.66 & 7.94 \\\\ \n \\hline \n $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\times$ & $\\surd$ & 0.67 & 6.08 \\\\\n \\hline \n $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\times$ & 0.69 & 6.81 \\\\ \n \\hline \n $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & $\\surd$ & 0.70 & 5.96 \\\\ \n \\hline \n \\end{tabularx}\n \\label{ablation}\n\\end{table*}\n\\subsection{Evaluation Metrics}\nMAE and MAPE are regarded as two evaluation criteria for training model. Here MAE~\\cite{willmott2005advantages} is defined as the Mean Absolute Error of corrected precipitation, and MAPE is a variant of MAE without rainless sampless ($<1mm$). In numerical weather prediction, threat score (TS) is a standard criterion for evaluating the accuracy of forecast~\\cite{mesinger2008bias} as follows:\n\\begin{equation} \\label{reconsLoss}\nTS_{\\rho} = H_{\\rho}\/(H_{\\rho} + M_{\\rho} +FA_{\\rho})\n\\end{equation}where $H_{.}$ is Hit (correction = 1, truth = 1), $M_{.}$ is Miss (correction = 0, truth = 1), and $FA_{.}$ is False Alarm (correction = 1, truth = 0), in which 1 is rainfall and 0 is rainless. Specifically, $\\rho$ is a threshold for splitting the range of Precipitation Intensity (PI) into two intervals and set $[0.1,1,10]$ for three different rainfall cases.\n\\subsection{Contrastive Experiments on ECbMi}\nWe list the assessment results of 8 methods and our models on the ECbMi show in Tab.~\\ref{contrastiveI}. The reported results are the average of 20 repetitions, each of which is the mean of predicted results on all batches. STAS outperforms all the other methods on five criteria and its $TS_{10}$ is $28.94\\%$ higher than the second highest result from OBA in this case. Meanwhile, the performance of FPN and LSTM can extract the spatial and temporal features severally beyond the traditional methods from third line to seventh line. Furthermore, FPN has preferable performance than LSTM. There are mainly two reasons for these phenomena above. 1) the performance of BCoP can be promoted either by learning temporal features or spatial features. 2) As for BCop, the spatial features are more important than temporal features. Furthermore, one possible reason is that FPN can predict rainfall utilized adaptive feature layer that has maximum likelihood~\\cite{lin2017feature}, but only learning temporal features in a fixed time scale for LSTM. Besides, the machine learning methods from the fourth to sixth line have somewhat better performance than original ECb forecasts because of utilizing more information from EC data.\n\\begin{table}[htp]\n\\caption{The 4 criteria between 3 machine learning methods and STAS on ECbT\/ECbM\/ECbH divided by precipitation intensity. N\/A(Not Applicable) is none of the samples in the current condition.}\n\\smallskip\n\\centering\n\\resizebox{150pt}{20mm}{\n\\smallskip\\begin{tabular}{c|c|c|c|c|c}\n\\hline \n &\\multicolumn{1}{|c}{Methods} &\\multicolumn{4}{|c}{Criteria} \\\\\n\\hline \nEcb & & MAPE & $TS_{0.1}$ & $TS_{1}$ & $TS_{10}$ \\\\\n\\hline \n\\multirow{4}*{EcbT} & OBA & 3.81 & 0.65 & N\/A & N\/A \\\\\n~ & LSTM & 3.79 & 0.55 & N\/A & N\/A \\\\\n~ & FPN & 2.14 & 0.60 & N\/A & N\/A \\\\\n~ & STAS & {\\bf 2.01} & 0.78 & N\/A & N\/A \\\\\n\\hline\n\\multirow{4}*{EcbM} & OBA & 8.39 & 0.53 & 0.49 & N\/A \\\\\n~ & LSTM & 8.01 & 0.51 & 0.47 & N\/A \\\\\n~ & FPN & 6.93 & 0.54 & 0.50 & N\/A \\\\\n~ & STAS & {\\bf 4.43} & 0.61 & {\\bf 0.59} & N\/A \\\\\n\\hline\n\\multirow{4}*{EcbH} & OBA & 13.44 & 0.21 & 0.21 & 0.09 \\\\\n~ & LSTM & 12.81 & 0.24 & 0.24 & 0.16 \\\\\n~ & FPN & 10.79 & 0.28 & 0.28 & 0.20 \\\\\n~ & STAS & {\\bf 7.05} & 0.38 & 0.38 & {\\bf 0.28} \\\\\n\\hline\n\\end{tabular}\n}\n\\label{contrastiveII}\n\\end{table}\n\\subsection{Contrastive Experiments on ECbT\/ECbM\/ECbH}\nFor investigating the influence of ST representation in different precipitation intensity, we compare the performance of the 3 targeted methods and STAS on ECBT, ECBM and ECbH as shown in Tab.~\\ref{contrastiveII}.\n\nWe prefer to select $MAPE$ instead of $MAE$ because correcting rainfall samples are our main purpose. Overall, The performance of all methods on $TS_{1}$ and $TS_{10}$ decreases, compared with the same $TS$s of these methods in Tab.~\\ref{contrastiveI}. The possible reason is that the large rainfall value is hard to correct because of the distribution of longspan and the few limited numbers of samples. Both FPN and LSTM have better performance in $TS_{10}$ than OBA because FPN can capture richer multi-scale spatial features and LSTM can obtain temporal dependency in EC. However, OBA only encodes deep representation. Furthermore, the $TS_{10}$ of LSTM on EcbH is $0.04$ lower than FPN in the same case. One possible explanation is that FPN can automatically select a scale layer with the largest confidence level as the predictive layer in the testing phase, whereas LSTM is a fixed ST scale before testing. Finally, It is worth noting that the $MAPE$s of all methods sharply rise with an increment of precipitation intensity. The reason is that the more samples heavy rainfall has, the bigger contribution MAPE according to its equation. Nevertheless, STAS nearly obtains all SOTA results on 3 Ecbs owing to the learning ability of optimized ST representation.\n\n\\subsection{Ablation Experiments}\nHere we perform ablation experiments to verify the effectiveness of each new-introduced components in STAS.\n\n\\noindent {\\bf Impact of SFM and TFM} \\quad The first two lines of Tab.~\\ref{ablation} show the effectiveness of SFM and TFM. It is obvious that the $TS_{1}$ sharply decreases after removing SFM or TFM, and the $TS_{1}$ of without ($w\/o$) SFM is $0.4$ lower than TFM in the same case. One possible reason is that the role of spatial scale is more important in BCoP than that of temporal scale. \n\n\\noindent {\\bf SFM vs TFM for ME modules}\\quad A similar conclusion can be obtained when we evaluate the SFM\/TFM impacts in 5ME modules if we see the results shown in the third line to the sixth line of Tab.~\\ref{contrastiveII}. Furthermore, it is obvious that the rainfall (R) module has a greater impact on $MAPE$ than other modules since it utilizes the historical observations of precipitation. \n \n\\noindent {\\bf Impact of deformable CNN and C3AE}\\quad The influences of Deformable CNN (D-CNN) and C3AE are shown in the seventh and eighth lines. MAE of $w\/o$ C3AE increases because C3AE is a regression method, quite like ordinal regression for solving the longspan range of precipitation distribution. Besides, the decrement of $TS_{1}$ on $w\/o$ D-CNN indicates that D-CNN does work in selecting optimized pixels in the process of learning spatial features. \n\n\n\\begin{figure}[t]\n\\centering\n\\includegraphics[width=\\columnwidth]{figure7.png} \n\\caption{ The visualized comparisons between predictions of precipitation and corresponding ground-truth in the same regions of specific stations. From left to right: the predicted precipitation on OBA and STAS respectively, and observed precipitation. The color-patch bar on the right of Figure is used for distinguishing precipitation intensity by changing light color to high color.}\n\\label{EC_vs}\n\\end{figure} \n\n\\subsection{Qualitative Analysis}\nWe visualize predictive results from 2 methods and corresponding observations in several examples shown in Fig.~\\ref{EC_vs}. Among the visualizations, the 2 red ovals on the left and middle pictures show that STAS can almost correct heavy rain (mazarine) accurately while OBA cannot correct exactly. Besides, the 2 orange ovals on the 2 same pictures above reveal that OBA has some errors on calibrating moderate rain (green), but STAS is successful for correcting it. All in all, STAS has better prediction performance than OBA in forecasting mesoscale or large-scale precipitation. \n\n\n\n\n\n\n\n\\iffalse\n\\noindent {\\bf Typhoon dataset} is the real-time observed typhoon data from MICAPS and records typhoon features such as maximum wind velocity and typhoon radius every $6h$. Specifically, we manually label these data according to maximum wind velocity~\\cite{cao2009wind}.\n\\fi\n\n\\iffalse\n\\noindent {\\bf Pre-processing} is divided into 3 procedures including feature selecting, slicing time-series ECMWF dataset in the given stations, and slicing the ECMWF data with typhoon flags. First, we utilize the Pearson Correlation to select the top 57 features that have the highest coefficients on BCoP from 601features in ECMWF. Second, we fix precipitation stations from MICAPS, and empirically slice the effective EC features around one station that are composed of the refined grid points in the geographical range of $\\mbox{[lat-2, lon-2]}$ to $\\mbox{[lat+2, lon+2]}$, then perform the time-slider in the time dimension to require time-series ECMWF dataset. Third, we select the ECMWF data that have the same time with typhoon and these ECMWF data are also included the $\\pm4$ range around the Lat-Lon of this typhoon radius. Meanwhile, we save the label and radius of the typhoon to guide the selection of ECMWF data and flag this typhoon data.\n\\fi\n\n\\iffalse\n\\begin{table}[htp]\n\\caption{The 5 criteria between 8 machine learning methods and STAS on ECbMi. $TS_{0.1}$ is TS score in $PI>0.1$. $TS_{1} | TS_{10}$ is $PI>1$ and $PI>10$ respectively. SVR is support vector regression, LR is linear regression, MLP is Multilayer Perceptron, LSTM is long short-term memory, and FCN is full convolutional network.}\n\\smallskip\n\\centering\n\\resizebox{240pt}{12mm}{\n\\smallskip\\begin{tabular}{c|c|c|c|c|c}\n\\hline \n &\\multicolumn{5}{|c}{Criterion($TS_{20\/30\/40}$regarding)} \\\\\n\\hline \nMethods& MAE & MAPE & $TS_{0.1}$ & $TS_{1}$ & $TS_{10}$ \\\\\n\\hline \nBCoP Benchmark Method~\\cite{ran2018evaluation} & 1.76$\\pm$0.13 & 3.18$\\pm$0.21 & 0.41$\\pm$0.04 & 0.35$\\pm$0.03 & 0.19$\\pm$0.04 \\\\\n\\hline\nSVR~\\cite{srivastava2015wrf} & 1.67$\\pm$0.12 & 3.04$\\pm$0.19 & 0.48$\\pm$0.04 & 0.41$\\pm$0.03 & 0.0$\\pm$0.0 \\\\\n\\hline\nLR~\\cite{hamill2012verification} & 1.73$\\pm$0.12 & 2.77$\\pm$0.15 & 0.35$\\pm$0.04 & 0.23$\\pm$0.03 & 0.21$\\pm$0.04 \\\\\n\\hline\nMLP~\\cite{yuan2007calibration} & 1.59$\\pm$0.11 & 2.54$\\pm$0.15 & 0.51$\\pm$0.04 & 0.41$\\pm$0.04 & 0.21$\\pm$0.04 \\\\\n\\hline\nLSTM~\\cite{yu2018coordinate} & (1.39$\\pm$0.08) & 2.37$\\pm$0.12 & 0.58$\\pm$0.04 & 0.48$\\pm$0.02 & 0.24$\\pm$0.02 \\\\\n\\hline\nFCN~\\cite{hu2018soil} & 1.26$\\pm$0.11 & 2.29$\\pm$0.14 & 0.61$\\pm$0.07 & 0.54$\\pm$0.04 & 0.27$\\pm$0.03 \\\\\n\\hline\nOBA\\iffalse~\\cite{hu2018soil}\\fi & 1.01$\\pm$0.05 & 2.21$\\pm$0.14 & 0.58$\\pm$0.03 & 0.53$\\pm$0.01 & 0.29$\\pm$0.01 \\\\\n\\hline\nSTAS(ours) & {\\bf 0.98$\\pm$0.03} & {\\bf 1.98$\\pm$0.10} & {\\bf 0.67$\\pm$0.03} & {\\bf 0.58$\\pm$0.02} & {\\bf 0.35$\\pm$0.02} \\\\\n\\hline\n\\end{tabular}\n}\n\\label{contrastive}\n\\end{table}\n\n\\fi\n\n\\iffalse The form of mean$\\pm$standard deviation. \\fi\n\\iffalse and changing curves on key hyperparameters respectively. \\fi\n\n\n\n\n\n\\section{Introduction}\n\nWeather forecast plays a crucial role in disaster monitoring and emergency disposal. Numerical Weather Prediction (NWP) based on the equations set of kinetic and thermodynamics~\\cite{liu2017kinetic} is often used to cope with sudden climate change and extreme weather beforehand. One of the representatives of progressive NWP in the global is the European Centre for medium-range weather forecasts (EC)~\\cite{ran2018evaluation}. However, the predictions of precipitation from EC suffer from some uncertainty intrinsic mechanisms of rainfall, e.g., the elusive physical process of rainfall. Therefore, Bias Correcting on Precipitation (BCoP) is the need to improve the forecast level of EC around a local area. \n\\begin{figure}[htbp]\n\\centering\n\\subfigure[$t-2$]{\n\\begin{minipage}[t]{0.3\\linewidth}\n\\centering\n\\includegraphics[width=1\\columnwidth,height=1.2\\textwidth]{figure1}\n\n\\end{minipage}%\n}%\n\\subfigure[$t-1$]{\n\\begin{minipage}[t]{0.3\\linewidth}\n\\centering\n\\includegraphics[width=1\\columnwidth,height=1.2\\textwidth]{figure2}\n\\end{minipage}%\n}%\n\\subfigure[$t$]{\n\\begin{minipage}[t]{0.3\\linewidth}\n\\centering\n\\includegraphics[width=1\\columnwidth,height=1.2\\textwidth]{figure3}\n\\end{minipage}\n}%\n\\centering\n\\caption{The visualization of EC precipitation of Eastern China in the 3 continuous timestamps with region segmentation, which can embody the spatiality and temporal granularity of precipitation.}\n\\label{ST}\n\\end{figure}\n\nClassical BCoP can be roughly divided into two categories: regression and parameter estimation. Regression methods can be regarded as a probabilistic model obeying a credible distribution judged by historical prior information such as expert experience~\\cite{hamill2008probabilistic}. And the parameter estimating methods figure out a set of optimal parameters in some functions such as Kalman filter~\\cite{monache2008kalman} via historical observation for better correcting. Note that in the era of restricted prior information or history observations, these two correcting methods less utilize the spatio-temporal weather regularity in EC~\\cite{hamill2012verification}. \n\nA feasible way to refine the performance of BCoP is to continuously learn the weather dynamic features from EC data and boosting the correcting ability with machine learning algorithms~\\cite{srivastava2015wrf}.\n\nNevertheless, shallow or low-level dynamic representation is not enough to significantly improve the performance of BCoP. Therefore, it is necessary to capture the high-level representation such as spatial and temporal-dependencies from EC data. To clarify this point, we visualize the colored EC precipitation region in the 3 continuous timestamps shown in Fig.~\\ref{ST}, and employ a Simple Linear Iterative Clustering (SLIC) is employed for segmenting different precipitation subregions by clustering the pixels in similar semantic information~\\cite{achanta2012slic}. From Fig.~\\ref{ST}, we can observe that pixels with the same color are naturally segmented into the same region, separated by the purple line, which reflects the spatiality of precipitation. Besides, an obvious movement for the positions of precipitation regions over time reflects varied temporal granularity. \n\nWe thus propose a deep ST Feature Auto-Selective (STAS) Model for learning ST representation. Further, we add two pre-trained modules termed Spatial Feature-selective Mechanism (SFM) and Temporal Feature-selective Mechanism (TFM) to STAS, in which five observations of meteorological elements (precipitation\/temperature\/pressure\/wind\/dew) are used to guide the selection of optimal ST scales shown in Fig.~\\ref{pipeline} for better extracting ST features. Besides, we integrate a binary classifier so that the precipitation prediction can be more precise. The final prediction is thus achieved by multiplying regression and classified results. The contributions of STAS are summarized as follows:\n\n\\begin{itemize}\n\\item {\\bf Spatial Adaptivity} The Spatial Feature-selective Mechanism (SFM) adaptively selects an optimal spatial scale of specific EC data for capturing the richer spatial representation, which can refine the performance of correcting, especially in the heavy rainfall, indicating its practicability in forecasting mesoscale or large-scale precipitation. \n\\item {\\bf Temporal Adaptivity} The Temporal Feature-selective Mechanism (TFM) is utilized for automatically choosing the optimal time-lagging sequence of time-series features of EC in line with the minimal loss value for acquiring the better temporal representation. \n\\item {\\bf Effective} Experiments indicate that our model achieved better prediction performance than the other 8 published methods on a BCOP benchmark dataset, especially for correcting the large-scale precipitation.\n\\end{itemize}\n\n\n\n\\section{STAS: A Spatio-Temporal Feature Auto-Selective Model}\n\nIn this section, we will introduce our proposed STAS for automatically selecting the spatial and temporal scales of meteorological features from European Centre for Medium-Range Weather Forecasts (EC) in detail. For better illustration, an overall pipeline of STAS is shown in Fig.~\\ref{pipeline}.\n\n\\begin{figure*}[!thbp] \n\\centering\n\\includegraphics[width=0.8\\textwidth,height=0.28\\textwidth]{figure4.png}\n\\caption{The framework is STAS. SFM and TFM is the mechanism of spatial and temporal feature-selective respectively. $s^{*}$ is adaptive spatial size of an EC data in one batch and $\\ell^{*}$ is adaptive time-lagging length of the encoded time-series features. $s^{*}$ and $\\ell^{*}$ are adaptively adjusted by SFM and TFM respectively. $\\mathcal{L}$s are main loss functions in STAS and RC is rainfall classifier. $\\mathcal{L}_{min}$s are minimal losses from the process of STAS.}\n\\label{pipeline}\n\\end{figure*}\n\n\\subsection{Notations and Methodology}\nFirst, we define several basic symbols for our method. Assuming that we have a total of $N$ surface observation stations from Eastern China. The EC data around one station can be set as the combination of the refined grid points in the geographical range from \\mbox{[$\\emph{da} - \\omega$, $\\emph{do} - \\omega$]} to \\mbox{[$\\emph{da} + \\omega$, $\\emph{do} + \\omega$]}, where \\emph{da} and \\emph{do} are the latitude and longitude of this station respectively, and $\\omega$ is the degree. Furthermore, the time-series EC data from the $i$th station at time level $t$ are defined as $(\\mat{X}^{i}_{t})_{\\ell} = [\\mat{X}_{t}^{u}, \\mat{X}_{t-1}, \\ldots, \\mat{X}_{t-\\tau}]_{\\ell}$, where $i \\in N$, and $\\ell$ is the length of time sequence and $\\ell = \\tau + 1$. $u$ is uniform spatial scale. In this study, the interval of the sequence $\\ell$ is 6h. Therefore, $(\\mat{X}^{i}_{t})_{\\ell}$ can be regarded as a four-dimensional tensor-form input including the dimensions of the channels of features, the length of the sequence, the height and width of features.\n\nWith these notations, we roughly build 3 subdivisions for performing Bias Correcting on Precipitation (BCoP) in order as follows:\n\\begin{eqnarray}\\label{eq_pipeline}\n \\hat{y}_{tp} &= & OR \\left \\{LSTM\\left ( E\\left ([\\mat{X}^{u}_{t}, \\ldots, \\mat{X}^{s^{*}_{\\tau}}_{t-\\tau}]\\right)_{\\ell^{*}}\\right) \\right\\} \\\\\n \\hat{y}_{rc} &= & RC(\\mat{X}^{u}_{t}) \\\\\n \\hat{y}_{t} &= & \\hat{y}_{tp} \\otimes \\hat{y}_{rc}\n\\end{eqnarray}\nwhere $s^{*}_{\\tau}$ denotes the optimized size (height $\\times$ width) of the EC data $\\mat{X}_{t-\\tau}$ through spatial feature-selective mechanism (SFM), which will be introduced in Sec.3.2. Parameter $\\ell^{*}$ denotes the refined time-lagging length for time-series features $[(\\mat{X}^{i}_{t})_{\\ell}]^{u}$ that are encoded (Sec.3.3) to a uniform ($u$) size via temporal feature-selective mechanism (TFM), which will be discussed in Sec.3.4. $E(\\cdot)$ is an encoder backbone and $LSTM(\\cdot)$ is a stacked ConvLSTM~\\cite{xingjian2015convolutional}. The acronym $OR(\\cdot)$ means an ordinal regression model~\\cite{zhu2018facial} is used for regressing corrected precipitation value in the end. Besides, we utilize the precipitation binary classifier $RC(\\cdot)$ for classifying raining or rainless samples. Finally, predicted result $\\hat{y}_{t}$ is obtained by multiplying the predicted precipitation value $\\hat{y}_{tp}$ and the classified result $\\hat{y}_{rc}$.\n\n\\subsection{Spatial Feature-Selective Mechanism}\n\\vspace{-0.1em}\nWhen lacking an instructive spatial scale, the EC features centered on all stations are empirically set to be a fixed spatial scale for prediction. However, these rules-of-thumb may impair predictive accuracy when there is a strong connection between features scale and precipitation intensity~\\cite{mu2003conditional}. Therefore, we propose a Spatial Feature-selective Mechanism (SFM) to adaptively search the optimal spatial scales of specific EC data based on observations of 5 Meteorological Elements (MEs) including precipitation, temperature, pressure, wind, dew. Concretely, we can find the optimal spatial scales by minimizing the total spatial losses in different scales, which are the summation of 5 spatial MSE losses between predictive MEs and their observations shown in Fig.~\\ref{frameSm}. The selection is formulated as:\n\\begin{eqnarray}\\label{spatialMe}\ns^{*} &= & \\mathop{\\arg\\min}_{s}\\mathcal{L}_{s} \\\\ \n\\mathcal{L}_{s} & = & \\sum_{i=0}^{n(MEs)}\\mathcal{L}_{MSE}(MSM_i(\\mat{X'}^{s}_{t}), y_{t}^{i})\n\\end{eqnarray}where $\\mathcal{L}_{s}$ is the spatial total loss in the scale $s$. $n(MEs)$ is the number of MEs. $MSM_i(\\cdot)$ is $i$-th ME Spatial Module and $\\mat{X'}^{s}_{t}$ is a given EC data that is scale $s$ in timestamp $t$. $y_{t}^{i}$ are the labels of $i$-th ME in timestamp $t$. Therefore, the goal is to search an optimal scale $s^{*}$ shown in Eq \\eqref{spatialMe}. Specifically, a deformable CNN layer is introduced for boosting the rep- resentational ability of 5 modules via learning the offsets of filters that are appropriate for capturing better features~\\cite{dai2017deformable}. \n\n\\subsection{Backbone with Denoising}\n\\vspace{-0.1em}\nThe EC features contain numerous noises~\\cite{xu2019towards} that can produce negative side effects for correcting. To solve this issue, we introduce a denoised Encoder-Decoder (E-D) as backbone shown in Fig.~\\ref{pipeline}. In the process of encoder, the Gaussian white noise is added to the encoded features that have uniform size $u$ via upsampling and downsampling. We also introduce a reconstruction loss utilized for calculating difference between encoded features and reconstructed features by decoder to optimize denoised ability of encoder as follows:\n\\begin{equation} \\label{reconsLoss}\n\\mat{W}_{E}^*, \\mat{W}_{D}^* = \\mathop{\\arg\\min}_{\\mat{W}_{E},\\mat{W}_{D}}{\\Vert E([\\mat{X}^{u}_{t}, \\ldots, \\mat{X}^{s^{*}_{\\tau}}_{t-\\tau}], \\mat{\\epsilon}) - D(Z_{t}) \\Vert}\n\\end{equation}where $\\mat{W}_{E}$ and $\\mat{W}_{D}$ are parameter matrices from encoder $E(\\cdot)$ and decoder $D(\\cdot)$ respectively. $[\\mat{X}^{u}_{t}, \\ldots, \\mat{X}^{s^{*}_{\\tau}}_{t-\\tau}]$ is time-series features with optimized spatial scale $\\vct{s^{*}}$ through SFM. $Z_{t}$ is hidden features by $E(\\cdot)$. Moreover, $\\mat{\\epsilon}$ is Gaussian noise.\n\\begin{figure}[t]\n\\centering\n\\includegraphics[width=1.4\\columnwidth, height=0.28\\textwidth]{figure5.png} \n\\caption{The structure of Spatial Feature-selective Mechanism (SFM), the visualizations of 5 Meteorological Elements (ME) are temperature, pressure, wind, dew and precipitation. $\\{ \\mathcal{L}_{s1}, \\ldots, \\mathcal{L}_{sw} \\}$ are MSE spatial total losses that are the sum ($\\oplus$) of 5 spatial MSE losses from ME spatial modules in different scales. SFM can select the adaptive spatial scale $s^{*}$ of specific EC data in $\\mat{X'}_{t}$. The deformable CNN is utilized for helping filters to operate the given pixels that can capture the better representation.}\n\\label{frameSm}\n\\end{figure} \n\\subsection{Temporal Feature-Selective Mechanism}\n\\vspace{-0.1em}\nIt is obvious that rainfall patterns in one station are not only related to EC features around this station, but also closely connected with historical features~\\cite{ciach2006analysis}. Therefore, the temporal Feature-selective Mechanism (TFM) is proposed for adaptively acquiring an optimized encoded features $[(\\mat{X}^{i}_{t})_{\\ell^*}]^{u} = [\\mat{X}^{u}_{t}, \\ldots, \\mat{X}^{u}_{t-\\tau}]_{\\ell^*}$, , from which more useful temporal representation can be learned. Formally, it is defined as: \n\\begin{eqnarray}\\label{temporalMe}\nl^{*} &= & \\mathop{\\arg\\min}_{\\ell}\\mathcal{L}_{T} \\\\\n\\mathcal{L}_{T} & = & \\sum_{i=0}^{n(MEs)}\\mathcal{L}_{MAE}(MTM_i([(\\mat{X}^{i}_{t})_{\\ell}]^{u}), y_{t}^{i})\n\\end{eqnarray}where $\\ell^{*}$ is the optimal time-lagging length and $y_{t}^{i}$ are labels of $i$-th MEs in timestamp $t$. $MTM_i(\\cdot)$ is $i$-th ME Temporal Module. Same as SFM, TFM can select optimal time-lagging sequence $\\ell^{*}$ by finding out minimum temporal total $MAE$ loss, which are the sum of 5 temporal $MAE$ losses between predictions of $MTM(\\cdot)$ and labels of MEs in different time-lagging sequences. The detailed structure of C3AE~\\cite{zhang2019c3ae} and 3DCNN~\\cite{zhang2017learning} are shown in Fig.~\\ref{frameTm}. In consideration of maldistribution of precipitation values~\\cite{xu2019towards}, an ordinal regression (OR) method~\\cite{zhu2018facial} is utilized for outputing regressing value of precipitation $\\hat{y}$ shown in Fig.~\\ref{pipeline}. OR may solve the problem of longspan range of precipitation values and convert a regression task into a multi-binary classification one to reduce the complexity of regression.\n\\begin{figure}[t]\n\\centering\n\\includegraphics[width=1.3\\columnwidth,height=0.28\\textwidth]{figure6.png} \n\\caption{ The structure of Temporal Feature-selective Mechanism (TFM). TFM can select adaptive time-lagging length $\\ell^*$ of one time-series features $(\\mat{X}^{i}_{t})_{\\ell}$. $\\{ \\mathcal{L}_{T1}, \\ldots, \\mathcal{L}_{Th} \\}$ are MAE temporal total losses that are the sum ($\\oplus$) of 5 MAE temporal losses from ME temporal modules in different time-lagging sequence. 3D CNN is utilized for capturing the patial-temporal dependency. Specifically, C3AE is a lightweight rank learning module and suited for regression distribution that has longspan range.}\n\\label{frameTm}\n\\end{figure} \n\n\\subsection{Training and Testing}\n{\\bf Training}\\quad In the training phase, we firstly pre-train the Spatial Feature-selective Module (SFM) and Temporal Feature-selective Module (TFM) so that these two modules can predict the MEs accurately. Secondly, we integrate SFM into our framework shown in Fig.~\\ref{pipeline} for predicting the MEs from historical EC data $\\mat{X'}_{t}^{s}$ and selecting the optimal scale $s^{*}$ that has minimal MSE loss. Similarly, we begin to train Temporal Feature-selective Module (TFM) and Encoder-Decoder (E-D) together, and use a specific ADAM optimizer for BP when adaptive time-lagging length $\\ell^{*}$ in one batch is selected by TFM. Meanwhile, the rainy classifier (RC) is cross-trained along with TFM\/SFM and E-D.\n\n\\noindent{\\bf Testing}\\quad In the testing phase, SFM plays a role in selecting the optimal scale $s^{*}$ and TFM selects the adaptive time-lagging length $\\ell^{*}$, then SFM is forward to calculate probabilities of all classifiers from ordinal regression. We select specific classifiers with large probabilities according to initialized interval $\\xi$ and transform these probabilities into regression value of precipitation $\\hat{y}_{tp}$ formulated as:\n\\begin{equation} \\label{temporalMe}\n\\hat{y}_{tp} = \\xi * \\sum\\limits_{v=0}^{c-1}(p_{v}\\ge\\xi)\n\\end{equation} where $p_{v}$ is classified probability of the $v$-th binary classifier. Besides, RC is forward to obtain the classified result $\\hat{y}_{rc}$. Final predicted result $\\hat{y}_{t}$ is required by multiplying SFM result $\\hat{y}_{tp}$ and RC result $\\hat{y}_{rc}$ as shown in the third equation of Eq \\eqref{eq_pipeline} in Sec.3.1.\n\n\n\n\n\n\\iffalse\nThe conditional nonlinear optimal perturbation (CNOP)~\\cite{mu2003conditional} is employed for searching a sensitive area in time dimension of history time level that can has an important effect of current prediction of precipitation. Some methods based on CNOP such as Ensemble Transforming Kalman Filter (ETKF) have some effect on finding sensitive area. However, these methods embrace excess intervention from experts and manual features, so increasing the overhead. \n\\fi \n\n\n\\iffalse\nTheoretically, the classified accuracy of RC is higher than OR because of only having one binary label. \n\\fi\n\n\\iffalse\n~\\cite{ciach2006analysis}\n\\fi\n\n\\iffalse\nSpecifically, the key role of RC is to enhance error-tolerant rate of SFM for improving performance of our model.\n\\fi\n\n\\section{Related Work}\n\n\\subsection{Bias Correcting on Precipitation} \nIn this section, we will give a brief survey on Bias Correcting on Precipitation (BCoP) and spatio-temporal (ST) pattern selection. Classical BCoP can be grouped into regression and parameter estimating methods. Regression methods~\\cite{hamill2008probabilistic} heavily depend on expert experience, which requires manually setting a threshold for the generation of probability, losing their flexibility and adaptivity. Meanwhile, the parameter estimating methods heuristically assess the key parameters set by trial and error in specific models. Nevertheless, traditional BCoP methods suffer from limited available priors and historical observations, and less utilize the clue of possible motion or dynamics existed in European Centre for medium-range weather forecasts (EC) data. To address these issues, some recently proposed models attempt to capture complex climate patterns from the large-scale EC data and optimize the model in line with observations~\\cite{srivastava2015wrf}. However, these methods neglect the potential dependency among variables in EC data, especially the ST dependency. \n\n\\subsection{Spatio-Temporal Pattern on EC Precipitation}\nExcept for qualitative analysis for reflecting the ST dependencies shown in Fig.~\\ref{ST}, it is proved that rainfall value in one location correlates with weather indicators such as temperature, pressure, wind and dew~\\cite{yapp1982model}. However, we do not know which scale around this location has a strong connection with the precipitation. Moreover, we know that the system on atmospheric dynamics is a spatio-temporal evolution system, in which physical field changes over time~\\cite{mu2003conditional}. Hence, capturing the ST representation would be a feasible way to improve the performance of BCoP. We consider employing an end-to-end deep-learning model based on multi-scale EC feature~\\cite{lin2017feature}, and adaptively select the features that have an optimal scale and take advantage of these features for correcting precipitation more accurately. To our knowledge, there is no report on how to adaptively extract the ST features based deep-learning model for BCoP in literature. This is the first time to adaptively extract ST representation end-to-end.\n \n\n\n\\iffalse \n\nIn deep learning domain, as we know, ST Long Short-Term Memory (ST-LSTM) is applied for action recognition by capturing the spatial features from joints of the trunk and temporal features from consecutive frames of action~\\cite{liu2016spatio}. Besides, a ST residual network is employed for traffic flow prediction~\\cite{zhang2017deep}.\n\nMoreover, conditional nonlinear optimal perturbation (CNOP), is utilized for identifying a sensitive time area that has a key influence on specific precipitation along the time dimension of this precipitation~\\cite{mu2003conditional}. \n\\fi \n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\\label{Sec. Introduction}\n\nIn his pioneering 1968 Ph.D. thesis, Leonard Parker discovered the surprising phenomenon that evolving cosmological spacetimes can produce quantum particles \\cite{Parker69, Parker71, ZeldovichStarobinsky77}. This work laid much of the theoretical framework for our current understanding of quantum fields in curved spacetimes (see, e.g., the textbooks \\cite{ParkerToms09, BirrellDavies82}). The phenomenon has since been investigated in astrophysical and cosmological contexts, leading to fundamental theoretical results including the emission of Hawking radiation by blackholes \\cite{Hawking75, Unruh76, Davies76, Page82, Brown86, Frolov87, Anderson93, Anderson94, Anderson95} and the generation of primordial fluctuations during inflation \\cite{Starobinsky79, Allen88, Sahni90, Mukhanov92, Souradeep92, Glenz09, Agullo11}.\n\nCosmological quantum particle production results from the shifting quantum vacuum state due to spacetime expansion. A given mode of a quantum field can begin in a state with no particles, but at a later time have a non-zero particle number expectation value. Parker showed that while this does not happen for massless conformally-coupled fields, it is generic for massive scalar fields of arbitrary coupling to the spacetime curvature. Typically, the particle production is significant when the particle mass $m$ is of the order of the expansion rate $H$, or when $m^{2} \\simeq \\dot{H}$. Simple dimensional arguments show that in a Universe with critical density $\\rho_c \\simeq 3H^{2}\/8\\pi G$, quantum particle production can contribute significantly to the energy density of the Universe at early epochs when $H$ is not too far below the Planck scale. \n\nThis effect naturally suggests that quantum particle production may have a significant impact on the expansion rate of the Universe at early times. The back-reaction problem consists in understanding how the energy density generated through this process alters the evolution of the background spacetime. This seemingly straightforward calculation is actually subtle due to two technical challenges of quantum fields in curved spacetimes. First, the definition of the quantum vacuum and particle are not formally well defined in time-varying spacetimes. Second, the energy-momentum tensor of a quantized field propagating on a curved background possesses a formally infinite expectation value, and must be regularized to yield physically sensible results. In homogeneous and isotropic spacetimes, both problems can be addressed via the method of adiabatic regularization \\cite{ParkerFulling74, FullingParker74, Bunch80, Fulling74, Christensen78, AndersonParker87}, which is particularly useful for numerical computations. The adiabatic notion of particle offers a clear way to track particle number in a spacetime which is expanding sufficiently slowly. This representation has the special property of defining a vacuum state which comes closest to the Minkowski vacuum when the background expansion is sufficiently slow. However, adiabatic regularization introduces several technical complications to the back-reaction problem. As a consequence, ambiguities arise in specifying initial conditions, and computationally the problem becomes very complicated with potential numerical instabilities \\cite{Anderson83, Anderson84, Anderson85, Anderson86, SuenAnderson87}. In practice, the ambiguities and complexities together have prevented any general numerical solution, although the test-field limit in which back-reaction effects are neglected has been investigated in several studies \\cite{Birrell78, Anderson00, Anderson05, Bates10, Habib99}.\n\nFormal developments clarified aspects of the adiabatic regularization approach, offering a clear separation between the energy density due to the field particle content and the divergent contributions from the zero-point energy \\cite{Habib99, AndersonMottola14, AndersonMottola14b, Anderson18}. However, the notion of adiabatic particle appearing in the energy density suffers from ambiguities. For a typical mode of a quantum field, its associated particle number during times of significant particle production depends on the perturbative order of adiabatic regularization employed. This is obviously an unphysical result, since the dominant term in the energy density is often just a simple function of the particle number density. Typically, successive orders of adiabatic regularization gives the particle number in a given mode as a divergent asymptotic series. Recent important papers by Dabrowski and Dunne \\cite{DabrowskiDunne14, DabrowskiDunne16} employed a remarkable result of asymptotic analysis \\cite{Dingle73, Berry90, Berry82, Berry89, Berry88} to sum the divergent series and provide a sensible notion of particle which is valid at all times. Since any remaining physical ambiguity is then removed from the problem, this result points the way to a general numerical solution.\n\nWe combine these recent results into an approach which numerically solves the quantum back-reaction problem in regimes with field energy density dominated by particle production. We then apply this technique to a toy cosmological model, that of a positive-curvature spacetime with a constant energy density (the closed de Sitter model). In the absence of any quantum fields or other particle content, this spacetime exhibits a bounce behavior, contracting to a minimum scale factor and then expanding again. Here we show explicitly that the existence of a massive scalar field in a particular mass range will cause large changes in the spacetime evolution: even if the contracting spacetime initially contains a quantum field in its adiabatic vacuum state, quantum particle production can create enough energy density to push the spacetime into a radiation crunch. Special values of the field mass can also delay but not eliminate the bounce. This appears to be the first general solution for the quantum back-reaction problem in cosmology.\n\nIn Section II we review standard results for quantized scalar fields propagating in cosmological spacetimes, while Section III discusses the adiabatic field representation and the semi-classical notion of adiabatic particle number. Section IV recasts quantum particle production in terms of the Stokes phenomenon of the complex-plane wave equation for specific modes, including interference between different modes. Section V formulates the back-reaction problem for scenarios in which the field particle content or particle production dominates the field energy density. Section VI outlines our numerical implementation of the mathematical results in Sections IV and V. Physical results for a closed de Sitter model are presented in Section VII. Finally, in Section VIII we discuss the prospects for more general situations, including quantum fields with spin and interacting quantum fields, and the possible relevance of quantum particle production to early-Universe models, including inflationary and bounce scenarios. Salient technical details are summarized in the Appendix. Natural units with $\\hbar = c = 1$ are adopted throughout.\n\n\\section{Scalar Fields in FLRW Spacetimes}\\label{Sec. QFT in FLRW}\n\nWe first summarize basic results for scalar fields in spatially isotropic and homogeneous spacetimes (see, e.g., \\cite{ParkerToms09, BirrellDavies82}). \nConsider a Universe described by the Friedmann-Lema\\^{i}tre-Robertson-Walker (FLRW) metric\n\\begin{equation}\\label{Eq. FLRW Metric}\n\t\\mathrm{d}s^{2} = g_{ab}\\mathrm{d}x^{a}\\mathrm{d}x^{b} = -\\mathrm{d}t^{2} + a^{2}(t) g_{ij}\\mathrm{d}x^{i}\\mathrm{d}x^{j} \\, ,\n\\end{equation}\nwith\n\\begin{equation}\\label{Eq. FLRW Spatial Metric}\n\tg_{ij}\\mathrm{d}x^{i}\\mathrm{d}x^{j} =\\frac{\\mathrm{d}r^{2}}{1 - K r^{2}} + r^{2} \\mathrm{d}\\theta^{2} + r^{2} \\sin^{2}\\theta \\, \\mathrm{d}\\varphi^{2} \\, .\n\\end{equation}\nHere $a(t)$ is the scale factor which describes the cosmological expansion history, and $K = -1, \\, 0, \\, +1$ corresponds to the curvature parameter of an open, flat, and closed Universe, respectively.\nThe non-vanishing components of the Ricci tensor $R_{ab}$ are\n\\begin{subequations}\\label{Eq. Ricci Tensor}\n\\begin{align}\n\tR_{00}(t) &= 3\\Big[ \\dot{H}(t) + H^{2}(t) \\Big] g_{00} \\, ,\\\\[8pt]\n\tR_{ij}(t) &= \\Big[ \\dot{H}(t) + 3 H^{2}(t) + 2K\/a^{2}(t) \\Big] g_{ij} \\, ,\n\\end{align}\n\\end{subequations}\nand the Ricci scalar $R = g^{ab}R_{ab}$ is\n\\begin{equation}\\label{Eq. Ricci Scalar}\n\tR(t) = 6 \\Big[ \\dot{H}(t) + 2 H^{2}(t) + K\/a^{2}(t) \\Big] \\, \n\\end{equation}\nwhere\n\\begin{equation}\\label{Eq. Hubble Parameter}\n\tH(t) \\equiv \\frac{\\dot{a}(t)}{a(t)} \n\\end{equation}\nis the Hubble parameter. Overdots indicate differentiation with respect to proper time $t$.\n\nWe are interested in the evolution of a free scalar field $\\Phi(t,\\, \\mathbf{x})$ of arbitrary mass and curvature coupling in this spacetime. The action for such a field can be expressed generically as\n\\begin{align}\\label{Eq. Field Action}\n\tS = -\\frac{1}{2} \\int \\sqrt{-g} \\, \\mathrm{d}^{4}x \\Big[ \\big( \\nabla_{a} \\Phi \\big) g^{ab} \\big( \\nabla_{b} \\Phi \\big) + \\Big. & \\\\\n\t \\Big. + m^{2} \\Phi^{2} + \\xi R \\Phi^{2} \\Big] & \\, , \\nonumber\n\\end{align}\nwhere $\\nabla_{a}$ is the covariant derivative, $g = \\det{(g_{ab})}$, $m$ is the field mass, and $\\xi$ is the field coupling to the spacetime curvature. Applying the variational principle to this action yields the equation of motion \n\\begin{equation}\\label{Eq. Field Equation of Motion}\n\t\\Big[ \\Box - m^{2} - \\xi R(t) \\Big]\\Phi(t,\\, \\mathbf{x}) = 0 \\, ,\n\\end{equation}\nwhere $\\Box = g^{ab}\\nabla_{a}\\nabla_{b}$ is the d'Alembert operator associated with the spacetime.\n\nDue to the homogeneity and isotropy of the background metric, the solutions of Eq.~(\\ref{Eq. Field Equation of Motion}) \ncan be separated into purely temporal and spatial parts. As a consequence, the quantized field operator can be written as\n\\begin{align}\\label{Eq. Field}\n\t\\hat{\\Phi}(t,\\, \\mathbf{x})=a^{-3\/2}(t) \\! \\int \\mathrm{d}\\mu(k) \\Big[ a_{\\mathbf{k}}^{\\phantom{\\dagger}}f_{k}(t)Y_{\\mathbf{k}}(\\mathbf{x}) + \\Big. & \\\\\n\t\\Big. + a_{\\mathbf{k}}^{\\dagger}f_{k}^{\\ast}(t)Y_{\\mathbf{k}}^{\\ast}(\\mathbf{x}) \\Big] & \\, , \\nonumber\n\\end{align}\nwhere the raising and lowering operators $a_{\\mathbf{k}}^{\\dagger}$ and $a_{\\mathbf{k}}^{\\phantom{\\dagger}}$ satisfy the canonical commutation relations\n\\begin{equation}\\label{Eq. Commutation Relations}\n\\left[ a_{\\mathbf{k}}^{\\phantom{\\dagger}} \\, , a_{\\mathbf{k^{\\prime}}}^{\\dagger} \\right] = \\delta_{\\mathbf{k},\\, \\mathbf{k^{\\prime}}} \\, ,\n\\end{equation}\nand $\\mathrm{d}\\mu(k)$ is a geometry-dependent integration measure given by \n\\begin{equation}\\label{Eq. Integration Measure}\n\t\\int \\mathrm{d}\\mu(k) =\n\t\\left\\{\\begin{aligned}\n\t\t\t&\\sum_{k=1}^{\\infty} k^2\\, , && \\text{for } K = +1 \\\\\n\t\t\t&\\int_{0}^{\\infty} k^2 \\,\\mathrm{d}k \\, , && \\text{for } K = 0,\\,-1 \\, .\n\t\t\\end{aligned}\\right.\n\\end{equation}\nThe functions $Y_\\mathbf{k}(\\mathbf{x})$ and $f_k(t)$ contain the spatial and temporal dependence of each $\\mathbf{k}$-mode. The harmonic functions $Y_\\mathbf{k}(\\mathbf{x})$ are eigenfunctions of the Laplace-Beltrami operator associated with the geometry of spatial hypersurfaces, while the mode functions $f_k(t)$ obey the harmnonic oscillator equation\n\\begin{equation}\\label{Eq. Mode Equation}\n\\ddot{f}_{k}(t) + \\Omega_{k}^{2}(t) f_{k}(t) = 0 \\, \n\\end{equation}\nwith the time-dependent frequency function \n\\begin{equation}\\label{Eq. Mode Frequency}\n\\Omega_{k}^{2}(t) = \\omega_{k}^{2}(t) + \\bigg( \\xi - \\frac{1}{6} \\bigg) R(t) - \\Bigg[ \\frac{\\dot{H}(t)}{2} + \\frac{H^{2}(t)}{4} \\Bigg] \\, ,\n\\end{equation}\nwhere\n\\begin{equation}\\label{Eq. Minkowski Frequency}\n\t \\omega_{k}(t) = \\Bigg[\\frac{k^{2}}{a^{2}(t)} + m^{2}\\Bigg]^{\\!1\/2} \\, .\n\\end{equation}\nThe complex mode functions $f_{k}(t)$ and $f^{\\ast}_{k}(t)$ also satisfy the Wronskian condition\n\\begin{equation}\\label{Eq. Wronskian Condition}\nf_{k}(t)\\dot{f}^{\\ast}_{k}(t) - \\dot{f}_{k}(t) f^{\\ast}_{k}(t) = i \\, .\n\\end{equation}\nIf Eq.~(\\ref{Eq. Wronskian Condition}) holds at some particular time $t$, then Eq.~(\\ref{Eq. Mode Equation}) guarantees it will also hold at all future\ntimes.\n\nThe quantization procedure outlined above naturally leads to the construction of the Fock space of field states. The base element of this space is the vacuum state, which is defined as the normalized state that is annihilated by all lowering operators:\n\\begin{equation}\\label{Eq. Vacuum}\na_{\\mathbf{k}}^{\\phantom{\\dagger}} \\left| 0 \\right\\rangle = 0 \\ \\ \\mathrm{and} \\ \\ \\left\\langle 0 | 0 \\right\\rangle = 1 \\, .\n\\end{equation}\nAll remaining states are generated from the vacuum by the successive application of raising operators, such as\n\\begin{equation}\n\\left| \\mathbf{k}_1 ,\\, \\mathbf{k}_2 ,\\, \\dots \\right\\rangle = a_{\\mathbf{k}_1}^{\\dagger} a_{\\mathbf{k}_2}^{\\dagger} \\dots \\left| 0 \\right\\rangle \\, ,\n\\end{equation}\nand normalized by the requirement of mutual orthonormality. In what follows, we will be interested in the family of field states which are spatially isotropic and homogeneous, as these constitute viable sources of the FLRW metric.\n\nThe field operator $\\hat{\\Phi}(t,\\, \\mathbf{x})$ admits numerous representations of the form shown in Eq.~(\\ref{Eq. Field}), each of which is associated with a different mode function pertaining to the set of solutions of Eq.~(\\ref{Eq. Mode Equation}). These representations are related: the complex mode functions $f_{k}(t)$ and $h_{k}(t)$ belonging to any two different representations can be expressed in terms of one another through the Bogolyubov transformations\n\\begin{subequations}\\label{Eq. Bogolyubov Modes}\n\\begin{align}\n\t{f}_{k}(t) &= \\alpha_{k}h_{k}(t) + \\beta_{k}h^{\\ast}_{k}(t) \\, ,\\\\[8pt]\n\t{f}_{k}^{\\ast}(t) &= \\beta_{k}^{\\ast}h_{k}(t) + \\alpha_{k}^{\\ast}h^{\\ast}_{k}(t) \\, ,\n\\end{align}\n\\end{subequations}\nwhere $\\alpha_{k}$ and $\\beta_{k}$ are known as Bogolyubov coefficients. Due to homogeneity and isotropy, these coefficients depend only on ${k = \\left| \\mathbf{k} \\right|}$. Substituting these expressions into Eq.~(\\ref{Eq. Field}) leads to similar transformations relating the raising and lowering operators belonging to these representations:\n\\begin{subequations}\\label{Eq. Bogolyubov Operators}\n\\begin{align}\na_{\\mathbf{k}}^{\\phantom{\\dagger}} &= \\alpha_{k}^{\\ast} b_{\\mathbf{k}}^{\\phantom{\\dagger}} - \\beta_{k}^{\\ast} b_{\\mathbf{k}}^\\dagger\\, , \\\\[8pt]\na_{\\mathbf{k}}^\\dagger &= \\alpha_{k} b_{\\mathbf{k}}^\\dagger - \\beta_{k} b_{\\mathbf{k}}^{\\phantom{\\dagger}} \\, ,\n\\end{align}\n\\end{subequations}\nfrom which it follows that the Bogolyubov coefficients must satisfy\n\\begin{equation}\\label{Eq. Bogolyubov Constraint}\n\t\\big| \\alpha_{k} \\big|^{2} - \\big| \\beta_{k} \\big|^{2} = 1 \\,\n\\end{equation}\nin order to guarantee that the commutation relations of Eq.~(\\ref{Eq. Commutation Relations}) are valid across all representations.\n\nA direct consequence of Eqs.~(\\ref{Eq. Bogolyubov Operators}) is that the notion of vacuum is not unique for a quantized field defined on a FLRW spacetime \\cite{BirrellDavies82, ParkerToms09}. This is evident from the following simple calculation, which shows that the vacuum defined in Eq.~(\\ref{Eq. Vacuum}) is not necessarily devoid of particles according to the number operator belonging to a different field representation:\n\\begin{align}\n\t\\mathcal{N}_{k} &= \\big\\langle 0 \\big| b_{\\mathbf{k}}^{\\dagger} b_{\\mathbf{k}}^{\\phantom{\\dagger}} \\big| 0 \\big\\rangle \\nonumber \\\\\n\t&= \\big| \\alpha_{k} \\big|^{2} \\big\\langle 0 \\big| a_{\\mathbf{k}}^{\\dagger} a_{\\mathbf{k}}^{\\phantom{\\dagger}} \\big| 0 \\big\\rangle + \\big| \\beta_{k} \\big|^{2} \\big\\langle 0 \\big| a_{\\mathbf{k}}^{\\phantom{\\dagger}} a_{\\mathbf{k}}^\\dagger \\big| 0 \\big\\rangle \\\\\n\t&= \\big| \\beta_{k} \\big|^{2} \\nonumber \\, .\n\\end{align}\nTherefore, different choices of representation inevitably lead to distinct notions of vacuum and, consequently, to distinct notions of particle. This result is a quite general feature of quantum field theory defined on curved spacetimes, and although it initially seems troublesome, it actually becomes useful in numerical back-reaction calculations. To that end, we introduce in the next section a particularly useful representation which defines the most physical notion of particle in a FLRW spacetime.\n\n\\section{Adiabatic Representation}\\label{Sec. Adiabatic Representation}\n\nDespite the multitude of available representations for a quantized scalar field defined on a FLRW spacetime, one particular choice referred to as the adiabatic representation stands out. This representation has the special property of defining a vacuum state which comes closest to the Minkowski vacuum when the background expansion is sufficiently slow. As a consequence, the adiabatic representation provides the most meaningful notion of physical particle in an expanding homogeneous and isotropic Universe. Here we discuss this representation \nclosely following Ref.~\\cite{Habib99}.\n\nThe adiabatic representation is characterized by mode functions which are the phase-integral solutions \\cite{Parker69, Parker71, ZeldovichStarobinsky77} of Eq.~(\\ref{Eq. Mode Equation}):\n\\begin{equation}\\label{Eq. Adiabatic Modes}\n\th_{k}(t) = \\frac{1}{\\sqrt{2 W_{k}(t)}} \\exp{\\!\\bigg( \\! -i \\! \\int^{t} W_{k}(s) \\, \\mathrm{d}s \\bigg)} \\, ,\n\\end{equation}\nwhere the integral in the exponent can be computed from any convenient reference time, and the function $W_{k}(t)$ is given by the formal asymptotic series\n\\begin{equation}\\label{Eq. Adiabatic W}\n\t W_{k}(t) \\equiv \\Omega_{k}(t) \\sum_{n=0}^{\\infty} \\varphi_{k,\\,2n}(t) \\, .\n\\end{equation}\nThe terms $\\varphi_{k,\\,2n}(t)$ are obtained by substituting Eqs.~(\\ref{Eq. Adiabatic W}) and (\\ref{Eq. Adiabatic Modes}) into Eq.~(\\ref{Eq. Mode Equation}). The expressions which ensue from these substitutions are standard results of the phase-integral method \\cite{FromanFroman96, FromanFroman02}; up to fourth order they are\n\\begin{subequations}\n\\begin{align}\n\t&\\varphi_{k,\\,0}(t) = 1 \\, , \\\\[8pt]\n\t&\\varphi_{k,\\,2}(t) = \\frac{1}{2} \\varepsilon_{k,\\,0}(t) \\, , \\\\[8pt]\n\t&\\varphi_{k,\\,4}(t) = -\\frac{1}{8}\\Big[ \\varepsilon_{k,\\,0}^{2}(t) + \\varepsilon_{k,\\,2}(t) \\Big] \\, ,\n\\end{align}\n\\end{subequations}\nfor which the quantities appearing on the right-hand sides are given by\n\\begin{subequations}\n\\begin{align}\n\t\\varepsilon_{k,\\,0}(t) &\\equiv \\Omega_{k}^{-3\/2}(t) \\, \\frac{d^{2}}{dt^{2}} \\! \\bigg[ \\Omega_{k}^{-1\/2}(t) \\bigg]\\, , \\\\[8pt]\n\t\\varepsilon_{k,\\,m}(t) &\\equiv \\bigg[ \\Omega_{k}^{-1}(t) \\, \\frac{d}{dt} \\bigg]^{\\!m} \\!\\!\\varepsilon_{k,\\,0}(t) \\, .\n\\end{align}\n\\end{subequations}\n\nIn a sense, the functions $W_{k}(t)$ capture the overall time dependence of each $\\mathbf{k}$-mode due the evolving FLRW metric, leaving behind only the Minkowski-like mode oscillations which take place on top of this background \\cite{DabrowskiDunne14, DabrowskiDunne16}. It is this property that makes the adiabatic mode functions $h_{k}(t)$ and $h^{\\ast}_{k}(t)$ such good templates for probing the particle content of fields evolving in cosmological spacetimes. This template role is made precise by the following time-dependent generalization of Eqs.~(\\ref{Eq. Bogolyubov Modes}) \\cite{Parker69, Parker71, ZeldovichStarobinsky77}, which expresses the field modes $f_{k}(t)$ as linear combinations of the adiabatic mode functions:\n\\begin{equation}\\label{Eq. Adiabatic Bogolyubov Mode}\n\t{f}_{k}(t) = \\alpha_{k}(t)h_{k}(t) + \\beta_{k}(t)h^{\\ast}_{k}(t) \\, ,\n\\end{equation}\nwhere the Bogolyubov coefficients $\\alpha_{k}(t)$ and $\\beta_{k}(t)$ are analogous to those appearing in Eqs.~(\\ref{Eq. Bogolyubov Modes})~and~(\\ref{Eq. Bogolyubov Operators}), but are here regarded as time-dependent quantities due to the fact that $h_{k}(t)$ and $h_{k}^{\\ast}(t)$ are merely approximate solutions of Eq.~(\\ref{Eq. Mode Equation}). In order to completely specify these coefficient functions, an additional expression must be provided. For that purpose, it is common to introduce a condition on the time derivative of the mode function which preserves the Wronskian relation of Eq.~(\\ref{Eq. Wronskian Condition}). In its most general form, this condition can be stated as \\cite{Habib99}\n\\begin{align}\\label{Eq. Adiabatic Bogolyubov Mode Derivative}\n\t\\dot{f}_{k}(t) = \\bigg[ -iW_{k}(t) + \\frac{V_{k}(t)}{2} \\bigg] \\alpha_{k}(t)h_{k}(t) \\,\\, & \\\\\n\t + \\, \\bigg[ iW_{k}(t) + \\frac{V_{k}(t)}{2} \\bigg] \\beta_{k}(t)h^{\\ast}_{k}(t)& \\, . \\nonumber\n\\end{align}\nHere the arbitrary function $V_{k}(t)$ contains the residual freedom in the definition of the adiabatic vacuum. In this work we will choose this function to be\n\\begin{equation}\\label{Eq. Adiabatic V}\n\tV_{k}(t) = -\\frac{\\dot{W}_{k}(t)}{W_{k}(t)} \\, ,\n\\end{equation}\nas this choice leads to important simplifications in the back-reaction problem. \n\nGathering Eqs.~(\\ref{Eq. Field}) and (\\ref{Eq. Adiabatic Bogolyubov Mode}), we find that the ladder operators associated with the adiabatic representation satisfy the transformations\n\\begin{subequations}\\label{Eq. Bogolyubov Adiabatic Operators}\n\\begin{align}\na_{\\mathbf{k}}^{\\phantom{\\dagger}} &= \\alpha_{k}^{\\ast}(t) b_{\\mathbf{k}}^{\\phantom{\\dagger}}(t) - \\beta_{k}^{\\ast}(t) b_{\\mathbf{k}}^{\\dagger}(t)\\, , \\\\[8pt]\na_{\\mathbf{k}}^\\dagger &= \\alpha_{k}(t) b_{\\mathbf{k}}^{\\dagger}(t) - \\beta_{k}(t) b_{\\mathbf{k}}^{\\phantom{\\dagger}}(t) \\, ,\n\\end{align}\n\\end{subequations}\nwhich in turn imply a time-dependent version of Eq.~(\\ref{Eq. Bogolyubov Constraint}), \n\\begin{equation}\\label{Eq. Adiabatic Bogolyubov Constraint}\n\t\\big| \\alpha_{k}(t) \\big|^{2} - \\big| \\beta_{k}(t) \\big|^{2} = 1 \\, .\n\\end{equation}\n\nFinally, it is useful to characterize field states according to the values of the non-trivial adiabatic bilinears $\\big\\langle b_{\\mathbf{k}}^{\\dagger}(t) b_{\\mathbf{k}}(t) \\big\\rangle$ and $\\big\\langle b_{\\mathbf{k}}(t) b_{\\mathbf{k}}(t) \\big\\rangle$. The first of these bilinears tracks the adiabatic particle content per comoving volume in the $\\mathbf{k}$-mode under consideration. Using the transformations established above by Eqs.~(\\ref{Eq. Bogolyubov Adiabatic Operators}), it follows that\n\\begin{align}\\label{Eq. Adiabatic Particles}\n\t\\mathcal{N}_{k}(t) &= \\big\\langle b_{\\mathbf{k}}^{\\dagger}(t) b_{\\mathbf{k}}(t) \\big\\rangle \\nonumber \\\\\n\t&= \\big| \\alpha_{k}(t) \\big|^{2} \\big\\langle a_{\\mathbf{k}}^{\\dagger} a_{\\mathbf{k}}^{\\phantom{\\dagger}} \\big\\rangle + \\big| \\beta_{k}(t) \\big|^{2} \\big\\langle a_{\\mathbf{k}}^{\\phantom{\\dagger}} a_{\\mathbf{k}}^\\dagger \\big\\rangle \\\\\n\t&= N_{k} + \\sigma_{k} \\, \\big| \\beta_{k}(t) \\big|^{2} \\nonumber \\, ,\n\\end{align}\nwhere $N_{k} = \\big\\langle a_{\\mathbf{k}}^{\\dagger} a_{\\mathbf{k}}^{\\phantom{\\dagger}} \\big\\rangle$ is a constant of motion which can be understood as the initial number of adiabatic particles per comoving volume populating the field mode of wavenumber $k$, and $\\sigma_{k} = 1 + 2N_{k}$ is the Bose-Einstein parameter responsible for stimulated particle production. The second bilinear can be expressed as\n\\begin{align}\\label{Eq. Adiabatic Interference}\n\t\\mathcal{M}_{k}(t) &= \\big\\langle b_{\\mathbf{k}}(t) b_{\\mathbf{k}}(t) \\big\\rangle \\nonumber \\\\\n\t&= \\alpha_{k}(t) \\beta^{\\ast}_{k}(t) \\big\\langle a_{\\mathbf{k}}^{\\dagger} a_{\\mathbf{k}}^{\\phantom{\\dagger}} \\big\\rangle + \\alpha_{k}(t) \\beta^{\\ast}_{k}(t) \\big\\langle a_{\\mathbf{k}}^{\\phantom{\\dagger}} a_{\\mathbf{k}}^\\dagger \\big\\rangle \\\\\n\t&= \\sigma_{k} \\, \\alpha_{k}(t) \\beta^{\\ast}_{k}(t) \\nonumber \\, .\n\\end{align}\n\nIn principle, these bilinears contain all the required information to track the field evolution and, consequently, the time dependence of the field energy density and pressure. In practice, however, these quantities suffer from an irreducible ambiguity which is particularly pronounced when $\\mathcal{N}_{k}(t)$ and $\\mathcal{M}_{k}(t)$ incur rapid changes, such as when particle production occurs. \nThe root of this issue can be traced back to the asymptotic representation of $W_{k}(t)$, which is usually handled by simply truncating the series in Eq.~(\\ref{Eq. Adiabatic W}) at a finite order. \nHowever, the values of the bilinears depend strongly on where the series is truncated if they are rapidly changing (see Ref.~\\cite{DabrowskiDunne14} for striking graphical representations).\nIn the next Section, we discuss a technique for finding the exact universal evolution for both adiabatic bilinears which the asymptotic series represents.\n\n\n\\section{Particle Production and the Stokes Phenomenon}\\label{Sec. Stokes}\n\nThe adiabatic representation introduced in the previous section provides an accurate description of the bilinears $\\mathcal{N}_{k}(t)$ and $\\mathcal{M}_{k}(t)$ whenever $\\left| \\varepsilon_{k,\\,0} \\right| \\ll 1$. The more severely this condition is violated, the more unreliable these adiabatic quantities become. Adiabatic particle production, for instance, coincides with the momentary violation of this condition, implying that the notion of particle remains uncertain until particle production ceases. Nonetheless, a universal notion of particle can be restored for all times when particle production events are understood in terms of the Stokes phenomenon.\n\nThe sharp transitions between asymptotic solutions of a given differential equation which are valid in different regions of the complex plane are termed the Stokes phenomenon. These regions are bounded by the so-called Stokes and anti-Stokes lines. In the context of a scalar field evolving in a FLRW spacetime, the differential equation of interest is the equation of motion for a given field mode extended to a complex time variable $z$:\n\\begin{equation}\\label{Eq. Complex Mode Equation}\n{f}^{\\prime\\prime}_{k}(z) + \\Omega_{k}^{2}(z) f_{k}(z) = 0 \\, ,\n\\end{equation}\nin which the primes stand for differentiation with respect to $z$, the proper time is given by $t \\equiv \\mathrm{Re}\\, z$, and $\\Omega_{k}(z)$ represents the analytic continuation of the time-dependent frequency of Eq.~(\\ref{Eq. Mode Frequency}). The Stokes lines associated with Eq.~(\\ref{Eq. Complex Mode Equation}) are those lines which emanate from the zeros (also known as turning points) and poles of $\\Omega_{k}(z)$ and along which $\\mathrm{Re}\\big[\\Omega_{k}\\mathrm{d}z\\big]~\\!\\!=~\\!\\!0$. An illustration of such a line is shown in Fig.~\\ref{Fig. Stokes Line}. The asymptotic solutions susceptible to the Stokes phenomenon are given by\n\\begin{equation}\\label{Eq. Phase-Integral Solution}\n\tf_{k}(z) = \\alpha_{k}(z) h_{k}(z) + \\beta_{k}(z) h^{\\ast}_{k}(z) \\, ,\n\\end{equation}\nwhere $h_{k}(z)$ and $h_{k}^{\\ast}(z)$ are the complex extensions of the adiabatic mode functions defined in the previous section. As this solution evolves across a Stokes line, the values of the Bogolyubov coefficients $\\alpha_{k}(z)$ and $\\beta_{k}(z)$ change abruptly. By Eqs.~(\\ref{Eq. Adiabatic Particles}) and (\\ref{Eq. Adiabatic Interference}), this implies a sudden change in the adiabatic bilinears and, in particular, the production of adiabatic particles. Remarkably, a result from asymptotic analysis guarantees the existence of a smooth universal form for this rapid transition between different asymptotic regimes. Below we outline the derivation of this important result and summarize the quantities which determine the functional form of such smooth Stokes jumps.\n\n\\begin{figure}[t!]\n \\includegraphics[width=0.49\\textwidth]{StokesLine.pdf}\n \\caption{A depiction of the Stokes line sourced by a conjugate pair of simple turning points $\\big( z_{0},\\, z_{0}^{\\ast}\\big)$ of a frequency function $\\Omega_{k}(z)$. Here the Stokes line crosses the real axis at the point $s_{0}$, which corresponds to the time at which particle production occurs for the mode of wavenumber $k$. The guiding lines show the directions for which the condition $\\mathrm{Re}\\big[\\Omega_{k}\\mathrm{d}z\\big] = 0$ is locally satisfied.}\n \\label{Fig. Stokes Line}\n\\end{figure}\n\nWe start by defining Dingle's singulant variable \\cite{Dingle73} anchored at $z_{0}$:\n\\begin{equation}\\label{Eq. Singulant}\nF_{k}^{(0)}(z) = 2i \\! \\int_{z_0}^{z} \\Omega_{k}(w) \\, \\mathrm{d}w \\, ,\n\\end{equation}\nwhere $z_0$ is a solution of $\\Omega_{k}(z) = 0$ which sources the Stokes line of interest, is closest to the real axis, and is located in the upper half-plane. The singulant is a convenient variable for tracking the change incurred by the Bogolyubov coefficients $\\alpha_{k}(z)$ and $\\beta_{k}(z)$ across a Stokes line. Indeed, it was shown by Berry \\cite{Berry90, Berry89, Berry88, Berry82} that these coefficients satisfy the following differential equations in the vicinity of a Stokes line:\n\\begin{subequations}\\label{Eq. Bogolyubov Couplers}\n\\begin{align}\n\t\\frac{d\\beta_{k}}{dF_{k}^{(0)}} &= C_{\\beta,\\,k}^{(0)}\\,\\alpha_{k}\\, , \\\\[8pt]\n\t\\frac{d\\alpha_{k}}{dF_{k}^{(0)}} &= C_{\\alpha,\\,k}^{(0)}\\,\\beta_{k} \\, ,\n\\end{align}\n\\end{subequations}\nwhere $C_{\\beta,\\,k}^{(0)}$ and $C_{\\alpha,\\,k}^{(0)}$ are coupling functions which depend on the order at which the series representation of $W_{k}(z)$ is truncated. A remarkable discovery by Dingle \\cite{Dingle73} states that the large-order terms in the asymptotic series of Eq.~(\\ref{Eq. Adiabatic W}) have a closed form given by\n\\begin{equation}\n\\varphi_{k,\\, 2n}(z) \\sim -\\frac{(2n - 1)!}{\\pi F_{k}^{(0)\\,2n}} \\quad \\text{for} \\quad n \\gg 1 \\, .\n\\end{equation}\nIt is clear from this result that the smallest term in such a series corresponds to $n \\approx \\big| F_{k}^{(0)} \\big|$. Terminating the series at this order leads to optimal closed form expressions for $C_{\\beta,\\,k}^{(0)}$ and $C_{\\alpha,\\,k}^{(0)}$, which can be substituted in Eqs.~(\\ref{Eq. Bogolyubov Couplers}) to yield the following universal behaviors for $\\beta_{k}(t)$ and $\\alpha_{k}(t)$ along the real axis and across the Stokes line under consideration:\n\\begin{subequations}\\label{Eq. Smooth Bogolyubov}\n\\begin{align}\n\t\\beta_{k}(t) &\\approx \\frac{i}{2} \\, \\mathrm{Erfc}\\Big( \\! -\\vartheta_{k}^{(0)}(t) \\Big) \\delta_{k}^{(0)}\\, , \\\\[8pt]\n\t\\alpha_{k}(t) &\\approx \\sqrt{1 + \\big| \\beta_{k}(t) \\big|^{2}} \\, ,\n\\end{align}\n\\end{subequations}\nwhere $\\vartheta_{k}^{(0)}(t)$ is a natural time evolution parameter which determines the sharpness of the Stokes jump, and $\\delta_{k}^{(0)}$ corresponds to the jump's amplitude. Both of these parameters are expressible in terms of the singulant variable evaluated over the real axis:\n\\begin{subequations}\n\\begin{align}\n\t&\\vartheta_{k}^{(0)}(t) = \\frac{\\phantom{\\sqrt{2 \\,\\,}}\\mathrm{Im} \\Big[ F_{k}^{(0)}(t) \\Big]}{\\sqrt{2 \\, \\mathrm{Re} \\Big[ F_{k}^{(0)}(t) \\Big]}}\\, , \\\\[8pt]\n\t&\\delta_{k}^{(0)} = \\exp\\Big( \\! -F_{k}^{(0)}(s_{0}) \\Big) \\, .\n\\end{align}\n\\end{subequations}\nHere $F_{k}^{(0)}(s_{0})$ is simply the singulant computed at the point $z=s_{0}$, where the Stokes line sourced by $z_{0}$ intersects the real axis, i.e., \n\\begin{equation}\\label{Eq. Amplitude}\nF_{k}^{(0)}(s_{0}) = 2i \\! \\int_{z_0}^{s_{0}} \\Omega_{k}(w) \\, \\mathrm{d}w = i \\! \\int_{z_0}^{z_0^{\\ast}} \\Omega_{k}(w)\\,\\mathrm{d}w \\, ,\n\\end{equation}\nwhere the last equality follows from the reality of $\\Omega_{k}(z)$ over the real axis. Putting together Eqs.~(\\ref{Eq. Adiabatic Particles}), (\\ref{Eq. Adiabatic Interference}), and (\\ref{Eq. Smooth Bogolyubov}) yields a universal functional form which describes the time evolution of the adiabatic bilinears associated with the field mode of wavenumber $k$:\n\\begin{subequations}\\label{Eq. Smooth Bilinears}\n\\begin{align}\n\t\\mathcal{N}_{k} &\\approx N_{k} + \\frac{\\sigma_{k}}{4} \\Big| \\mathrm{Erfc}\\Big( \\! -\\vartheta_{k}^{(0)} \\Big) \\delta_{k}^{(0)} \\Big|^{2}\\, , \\\\[8pt]\n\t\\mathcal{M}_{k} &\\approx -i \\, \\frac{\\sigma_{k}}{2} \\Big[ \\mathrm{Erfc}\\Big( \\! -\\vartheta_{k}^{(0)} \\Big) \\delta_{k}^{(0)} \\Big] \\Big[ 1 + \\mathcal{N}_{k} \\Big]^{\\!1\/2}\n\\, .\n\\end{align}\n\\end{subequations}\n\nThese results can be further generalized to account for multiple Stokes line crossings, as well as the interference effects between them \\cite{DabrowskiDunne14}. Define the accumulated phase between the first and the $p$-th pair of zeros of $\\Omega_{k}(z)$ as\n\\begin{equation}\\label{Eq. Phase}\n\\theta_{k}^{(p)} = \\int_{s_0}^{s_p} \\Omega_{k}(w)\\,\\mathrm{d}w \n\\end{equation}\nwhere $s_p$ corresponds to the point where the Stokes line associated with the $p$-th conjugate pair of zeros crosses the real axis. The functions which describe both adiabatic bilinears are then given by\n\\begin{subequations}\\label{Eq. Bilinears Multiple Crossings}\n\\begin{align}\n\t\\mathcal{N}_{k} &\\approx N_{k} + \\frac{\\sigma_{k}}{4} \\bigg| \\sum_{p} \\mathrm{Erfc}\\Big( \\! -\\vartheta_{k}^{(p)} \\Big) \\delta_{k}^{(p)} \\exp\\Big(2i\\theta_{k}^{(p)} \\Big) \\bigg|^{2}\\, , \\\\[8pt]\n\t\\mathcal{M}_{k} &\\approx -i \\, \\frac{\\sigma_{k}}{2} \\bigg[ \\sum_{p} \\mathrm{Erfc}\\Big( \\! -\\vartheta_{k}^{(p)} \\Big) \\delta_{k}^{(p)} \\exp\\Big( \\! -2i\\theta_{k}^{(p)} \\Big) \\bigg] \\times \\nonumber \\\\\n\t & \\,\\,\\;\\; \\quad \\quad \\times \\! \\bigg[ 1 + \\mathcal{N}_{k} \\bigg]^{\\!1\/2}\n\\end{align}\n\\end{subequations}\nwith $\\delta_{k}^{(p)}$ and $\\vartheta_{k}^{(p)}(t)$ being the amplitude and time evolution parameter associated with the $p$-th Stokes line.\n\nTherefore, by monitoring the turning points and Stokes lines which accompany each mode's frequency function on the complex plane, we can track the evolution of the adiabatic bilinears related to any physically acceptable field state. In the next section we examine how this evolution affects the Universe's scale factor through the semi-classical Einstein equations.\n\n\\onecolumngrid\n\\section{The Semi-Classical Einstein Equations}\\label{Sec. Semi-Classical Einstein}\n\nIf cosmological quantum particle production occurs at a sufficiently high rate, it can in principle back-react on the cosmic evolution through the semi-classical Einstein equations\n\\begin{equation}\\label{Eq. Einstein Equations}\n\tR_{ab} - \\frac{1}{2}R \\, g_{ab} + \\Lambda \\, g_{ab} = M^{-2} \\big\\langle \\hat{T}_{ab}\\big\\rangle \n\\end{equation}\nwhere $\\Lambda$ represents the cosmological constant, ${M = ({8 \\pi G})^{-1\/2}}$ stands for the reduced Planck mass, and $\\big\\langle \\hat{T}_{ab}\\big\\rangle$ corresponds to the expectation value of the energy-momentum tensor operator, including contributions both from the scalar field we are considering plus any other stress-energy sources. The canonical expression for $\\hat{T}_{ab}$ due to the scalar field is constructed by varying the action in Eq.~(\\ref{Eq. Field Action}) with respect to the metric $g_{ab}$, and subsequently substituting the field operator $\\hat{\\Phi}$ from Eq.~(\\ref{Eq. Field}) into the resulting expression:\n\\begin{equation}\\label{Eq. Energy-Momentum Tensor}\n\t\\hat{T}_{ab} = \\big( \\nabla_{a} \\hat{\\Phi} \\big) \\big( \\nabla_{b} \\hat{\\Phi} \\big) - \\frac{1}{2}g_{ab} \\big( \\nabla^{c} \\hat{\\Phi} \\big) \\big( \\nabla_{c} \\hat{\\Phi} \\big) + \\xi \\bigg[ g_{ab} \\, \\Box - \\nabla_{a}\\nabla_{b} + R_{ab} - \\frac{1}{2}R \\, g_{ab} - \\frac{m^{2}}{2}g_{ab} \\bigg] \\hat{\\Phi}^{2} \\, .\n\\end{equation}\nFor a statistically homogeneous and isotropic field state, $\\big\\langle \\hat{T}_{ab}\\big\\rangle$ is equivalent to the energy-momentum tensor of a perfect fluid for which the field energy density and pressure are given, respectively, by $\\rho(t) = \\big\\langle \\hat{T}_{00} \\big\\rangle$ and $P(t) = \\frac{1}{3}g^{ij}\\big\\langle \\hat{T}_{ij} \\big\\rangle$. As a consequence, such a state naturally sources an FLRW metric, reducing Eq.~(\\ref{Eq. Einstein Equations}) to the usual Friedmann equations:\n\\begin{subequations}\\label{Eq. Friedmann Equations}\n\\begin{align}\n\t& H^{2}(t) = \\frac{1}{3} M^{-2} \\rho(t) + \\frac{\\Lambda}{3} - \\frac{K}{a^{2}(t)} \\label{Eq. First Friedmann Equation} \\\\[8pt]\n\t&\\dot{H}(t) + H^{2}(t) = -\\frac{1}{6} M^{-2} \\Big[ \\rho(t) + 3P(t) \\Big] + \\frac{\\Lambda}{3} \\, .\n\\end{align}\n\\end{subequations}\nFurthermore, it can be shown that $\\big\\langle \\hat{T}_{ab}\\big\\rangle$ is covariantly conserved, resulting in the cosmological continuity equation\n\\begin{equation}\\label{Eq. Cosmological Continuity}\n\t\\dot{\\rho}(t) + 3H(t)\\Big[ \\rho(t) + P(t) \\Big] = 0 \\, .\n\\end{equation}\n\nHowever, at this stage these equations are merely formal, because both the pressure and energy density of a quantized field are in general divergent and need to be regularized. \nIn a FLRW spacetime, these divergencies can be partially isolated by expressing $\\rho(t)$ and $P(t)$ in terms of the adiabatic bilinears $\\mathcal{N}_{k}(t)$ and $\\mathcal{M}_{k}(t)$ \\cite{Habib99}. Substituting Eqs.~(\\ref{Eq. Field}), (\\ref{Eq. Adiabatic Bogolyubov Mode}), and (\\ref{Eq. Adiabatic Bogolyubov Mode Derivative}) into the expectation value of Eq.~(\\ref{Eq. Energy-Momentum Tensor}) and collecting terms with the same adiabatic factor gives\n\\begin{subequations}\\label{Eq. Energy Density and Pressure}\n\\begin{align}\n\t&\\rho(t) = \\big\\langle \\hat{T}_{00} \\big\\rangle = \\frac{1}{4 \\pi a^{3}(t)} \\! \\int d\\mu(k) \\Bigg\\{ \\rho^{\\mathcal{N}}_{k}\\!(t) \\bigg[ \\mathcal{N}_{k}(t) + \\frac{1}{2} \\bigg] + \\rho^{\\mathcal{R}}_{k}(t)\\,\\mathcal{R}_{k}(t) + \\rho^{\\mathcal{I}}_{k}(t) \\, \\mathcal{I}_{k}(t) \\Bigg\\} \\label{Eq. Energy Density} \\\\[8pt] \n\t& P(t) = \\frac{1}{3} g^{ij} \\big\\langle \\hat{T}_{ij} \\big\\rangle = \\frac{1}{4 \\pi a^{3}(t)} \\! \\int d\\mu(k) \\Bigg\\{ P^{\\mathcal{N}}_{k}\\!(t) \\bigg[ \\mathcal{N}_{k}(t) + \\frac{1}{2} \\bigg] + P^{\\mathcal{R}}_{k}(t)\\,\\mathcal{R}_{k}(t) + P^{\\mathcal{I}}_{k}(t) \\, \\mathcal{I}_{k}(t) \\Bigg\\} \\label{Eq. Pressure} \\, .\n\\end{align}\n\\end{subequations}\nThe terms proportional to $\\mathcal{N}_{k}(t)$ capture the contribution to the energy density and pressure due to the evolving distribution of adiabatic particles populating the field modes, while the quantum interference terms contain ${\\mathcal{R}_{k}(t) = \\mathrm{Re}\\big[ \\mathcal{M}_{k}(t) \\big]}$ and ${\\mathcal{I}_{k}(t) = \\mathrm{Im}\\big[ \\mathcal{M}_{k}(t) \\big]}$ \\footnote{Here the definitions for $\\mathcal{R}_{k}(t)$ and $\\mathcal{I}_{k}(t)$ might seem to differ from those found in the literature by a phase factor, but this factor is implicit in our definitions for $\\alpha_{k}(t)$ and $\\beta_{k}(t)$ obtained from asymptotic analysis}. The prefactors in each of these terms are defined as\n\\begin{subequations}\\label{Eq. Adiabatic Factors}\n\\begin{align}\n\t&\\rho^{\\mathcal{N}}_{k}\\!(t) \\equiv \\frac{1}{W_{k}(t)} \\Bigg\\{ W^{2}_{k}(t) + \\omega^{2}_{k}(t) + \\frac{1}{4} \\Big[ V_{k}(t) - H(t) \\Big]^{\\!2} + \\big(6\\xi - 1\\big) \\bigg[ H(t) V_{k}(t) - 2H^{2}(t) + \\frac{K}{a^{2}(t)} \\bigg] \\Bigg\\}\\, , \\\\[8pt]\n\t&\\rho^{\\mathcal{R}}_{k}(t) \\equiv \\frac{1}{W_{k}(t)} \\Bigg\\{ - W^{2}_{k}(t) + \\omega^{2}_{k}(t) + \\frac{1}{4} \\Big[ V_{k}(t) - H(t) \\Big]^{\\!2} + \\big(6\\xi - 1\\big) \\bigg[ H(t) V_{k}(t) - 2H^{2}(t) + \\frac{K}{a^{2}(t)} \\bigg] \\Bigg\\}\\, , \\\\[8pt]\n\t&\\rho^{\\mathcal{I}}_{k}(t) \\equiv V_{k}(t) - H(t) +2H(t)\\big(6\\xi - 1\\big)\\, , \\\\[8pt]\n\t&P^{\\mathcal{N}}_{k}\\!(t) \\equiv \\frac{1}{3W_{k}(t)} \\Bigg\\{ W^{2}_{k}(t) + \\omega^{2}_{k}(t) - 2m^{2} + \\frac{1}{4} \\Big[ V_{k}(t) - H(t) \\Big]^{\\!2} + \\frac{1}{3}\\big(6\\xi - 1\\big)^{\\!2}R(t) + \\Bigg. \\, \\\\\n\t & \\quad \\quad \\quad \\quad \\quad \\quad \\quad \\quad \\Bigg. + \\big(6\\xi - 1\\big) \\bigg[ -2W^{2}_{k}(t) - \\frac{1}{2}V^{2}_{k}(t) + 4H(t)V_{k}(t) + 2\\omega^{2}_{k}(t) + 2\\dot{H}(t) + \\frac{K}{a^{2}(t)} - \\frac{5}{2}H^{2}(t) \\bigg] \\Bigg\\}\\, , \\nonumber \\\\[8pt]\n\t&P^{\\mathcal{R}}_{k}(t) \\equiv \\frac{1}{3W_{k}(t)} \\Bigg\\{ -W^{2}_{k}(t) + \\omega^{2}_{k}(t) - 2m^{2} + \\frac{1}{4} \\Big[ V_{k}(t) - H(t) \\Big]^{\\!2} + \\frac{1}{3}\\big(6\\xi - 1\\big)^{\\!2}R(t) + \\Bigg.\\, \\\\\n\t & \\quad \\quad \\quad \\quad \\quad \\quad \\quad \\quad \\Bigg. + \\big(6\\xi - 1\\big) \\bigg[ 2W^{2}_{k}(t) - \\frac{1}{2}V^{2}_{k}(t) + 4H(t)V_{k}(t) + 2\\omega^{2}_{k}(t) + 2\\dot{H}(t) + \\frac{K}{a^{2}(t)} - \\frac{5}{2}H^{2}(t) \\bigg] \\Bigg\\}\\, , \\nonumber \\\\[8pt]\n\t&P^{\\mathcal{I}}_{k}(t) \\equiv \\frac{1}{3}\\Big[ V_{k}(t) - H(t) \\Big] + \\frac{2}{3} \\big(6\\xi - 1\\big) \\Big[ 4H(t) - V_{k}(t) \\Big] \\, .\n\\end{align}\n\\end{subequations}\nFor adiabatic field states, the contributions to the energy density and pressure due to the real bilinears $\\mathcal{N}_{k}(t)$, $\\mathcal{R}_{k}(t)$, and $\\mathcal{I}_{k}(t)$ are always finite. This implies that the divergencies in Eqs.~(\\ref{Eq. Energy Density}) and (\\ref{Eq. Pressure}) are isolated in the vacuum-like terms characterized by the $\\frac{1}{2}$ factors, henceforth identified as\n\\begin{subequations}\\label{Eq. Energy and Pressure Vacuum}\n\\begin{align}\n\t\\rho_{\\mathrm{vac}}(t) &\\equiv \\frac{1}{8 \\pi a^{3}(t)} \\! \\int d\\mu(k)\\,\\rho^{\\mathcal{N}}_{k}\\!(t) \\label{Eq. Vacuum Energy} \\\\[8pt]\n\tP_{\\mathrm{vac}}(t) &\\equiv \\frac{1}{8 \\pi a^{3}(t)} \\! \\int d\\mu(k)\\,P^{\\mathcal{N}}_{k}\\!(t) \\label{Eq. Vacuum Pressure} \\, .\n\\end{align}\n\\end{subequations}\n\nRegularization consists precisely in controlling the divergent behavior of $\\rho_{\\mathrm{vac}}(t)$ and $P_{\\mathrm{vac}}(t)$ so as to obtain finite expressions for $\\rho(t)$ and $P(t)$ which still satisfy the cosmological continuity equation. Adiabatic regularization achieves this result by subtracting the fourth-order phase-integral expansions of $\\rho^{\\mathcal{N}}_{k}\\!(t)$ and $P^{\\mathcal{N}}_{k}\\!(t)$ from the integrands of Eqs.~(\\ref{Eq. Vacuum Energy}) and (\\ref{Eq. Vacuum Pressure}), respectively \\cite{ParkerFulling74, AndersonParker87}. From a technical perspective, however, this procedure introduces significant challenges to the numerical implementation of the semi-classical Friedmann equations. Chief among these is the appearance of higher-order time derivatives of $H(t)$ in the integrands of Eqs.~(\\ref{Eq. Energy Density and Pressure}), turning the semi-classical Friedmann equations into a system of integro-differential equations which is not amenable to standard numerical treatments. We circumvent this difficulty by employing an alternative regularization scheme which, albeit cruder, yields a good approximation to $\\rho(t)$ and $P(t)$ in regimes dominated by particle production.\n\nCentral to the regularization approach adopted here is the realization that $\\rho_{\\mathrm{vac}}(t)$ and $P_{\\mathrm{vac}}(t)$ independently satisfy the cosmological continuity equation as long as the function $V_{k}(t)$ has the form established in Eq.~(\\ref{Eq. Adiabatic V}) (see the Appendix). It follows that the vacuum contributions to $\\rho(t)$ and $P(t)$ can be discarded in their entirety while still ensuring that Eq.~(\\ref{Eq. Cosmological Continuity}) remains valid. Despite its simplicity, this procedure yields a good approximation to the field energy density and pressure provided the adiabatically regularized integrands of Eqs.~(\\ref{Eq. Energy Density and Pressure}) are dominated by the real adiabatic bilinears $\\mathcal{N}_{k}(t)$, $\\mathcal{R}_{k}(t)$, and $\\mathcal{I}_{k}(t)$. Therefore, in what follows we take the regularized expressions for the energy density and pressure to be\n\\begin{subequations}\\label{Eq. Reg. Energy Density and Pressure}\n\\begin{align}\n\t\\rho(t) &\\approx \\frac{1}{4 \\pi a^{3}(t)} \\! \\int d\\mu(k) \\Bigg\\{ \\rho^{\\mathcal{N}}_{k}\\!(t) \\, \\mathcal{N}_{k}(t) + \\rho^{\\mathcal{R}}_{k}(t)\\,\\mathcal{R}_{k}(t) + \\rho^{\\mathcal{I}}_{k}(t) \\, \\mathcal{I}_{k}(t) \\Bigg\\} \\label{Eq. Reg. Energy Density} \\\\[8pt] \n\tP(t) &\\approx \\frac{1}{4 \\pi a^{3}(t)} \\! \\int d\\mu(k) \\Bigg\\{ P^{\\mathcal{N}}_{k}\\!(t) \\, \\mathcal{N}_{k}(t) + P^{\\mathcal{R}}_{k}(t)\\,\\mathcal{R}_{k}(t) + P^{\\mathcal{I}}_{k}(t) \\, \\mathcal{I}_{k}(t) \\Bigg\\} \\label{Eq. Reg. Pressure} \\, ,\n\\end{align}\n\\end{subequations}\nwhere the factors $\\rho^{\\mathcal{N}}_{k}\\!(t)$, $\\rho^{\\mathcal{R}}_{k}(t)$, and $\\rho^{\\mathcal{I}}_{k}(t)$, $P^{\\mathcal{N}}_{k}\\!(t)$, $P^{\\mathcal{R}}_{k}(t)$, and $P^{\\mathcal{I}}_{k}(t)$ are computed by truncating the asymptotic series Eq.~(\\ref{Eq. Adiabatic W}) for $W_{k}(t)$ at its optimal order.\n\nFinally, the regularization of $\\big\\langle \\hat{T}_{ab}\\big\\rangle$ also induces the renormalization of the gravitational coupling constants $G$ and $\\Lambda$. Moreover, self-consistency demands the introduction of a covariantly conserved tensor composed of fourth-order derivatives of the metric into the semi-classical Einstein equations \\cite{FullingParker74, Bunch80}. This tensor is accompanied by a new unknown coupling constant whose renormalization assimilates the ultra-violet divergence in the field energy-momentum tensor. For simplicity, in this work we assume this new coupling constant to be renormalized to zero, thus preserving the form of Eq.~(\\ref{Eq. Einstein Equations}). Non-zero values for this coupling constant will be considered elsewhere.\n\nTaken together, Eqs.~(\\ref{Eq. Bilinears Multiple Crossings}), (\\ref{Eq. Friedmann Equations}), and (\\ref{Eq. Reg. Energy Density and Pressure}) describe the coupled field evolution and cosmic evolution in regimes dominated by particle production. In the next section we present an algorithm which numerically solves this system of equations.\\\\\n\\twocolumngrid\n\n\\section{Numerical Implementation}\\label{Sec. Numerical Implementation}\n\nThe semi-classical Friedmann equations can be formulated as a discretized initial value problem. We take the domain of numerical integration to be a band of the complex plane which is bisected by the real $t$ axis. As illustrated in Figure \\ref{Fig. Numerical Analytic Continuation}, this band is discretized by a uniformly spaced grid where the real-valued entries $t_{j}$ label the physical time. Initial conditions are set by an appropriately chosen functional form for the scale factor $a(t)$ which not only admits an adiabatic field state at the initial time $t_{0}$, but which is also consistent with our choice for the initial distribution of adiabatic particles $\\mathcal{N}_{k}(t_{0}) = N_{k}$ populating the field modes. In addition, we require that\n\\begin{equation}\\label{Eq. Particle Distribution Constraint}\n\tN_{k} < \\mathcal{O}(k^{-3}) \\quad \\text{as} \\quad k \\rightarrow \\infty\n\\end{equation}\nin order to ensure that both the energy density and pressure associated with the initial particle distribution are finite. \n\nWe use a standard finite-difference scheme to step $a(t)$, $H(t)$, and $\\dot{H}(t)$ along the real axis, and employ B-splines to scan the Stokes geometry on the complex plane. The latter is accomplished by generating a numerical sample of $\\Omega_{k}(t)$ through Eq.~(\\ref{Eq. Mode Frequency}), and subsequently performing high-order B-spline interpolations to construct a truncated Taylor polynomial for this function over the real line up to the value of $t$ in the current time step. Due to the analyticity of $\\Omega_{k}(t)$, this series representation is also valid on the complex plane, and thus encodes the analytical continuation of the frequency function. Explicitly, given a grid point $z_{ij}$ on the discretized plane, we compute $\\Omega_{k}(z_{ij})$ through the expression\n\\begin{equation}\\label{Eq. Taylor}\n\\Omega_{k}(z_{ij}) \\approx \\sum_{n = 0}^{T} \\frac{1}{{n!}} \\big(z_{ij} - t_{j} \\big)^{n} \\, \\Omega_{k}^{(n)}(t_{j}) \\, ,\n\\end{equation}\nwhere $t_{j} = \\mathrm{Re}\\,z_{ij}$, as depicted in Figure \\ref{Fig. Numerical Analytic Continuation}. The numerical derivatives $\\Omega_{k}^{(n)}$ are extracted from B-spline interpolations over the real axis, and $T$ corresponds to a truncation order which depends on the density of grid points lying over the real axis. In addition, we feed Eq.~(\\ref{Eq. Taylor}) to a Pad\\'{e} approximant \\cite{Bender13} routine to accelerate its convergence and improve its accuracy. Once this approximate representation of the frequency function has been computed over the discretized plane, it can be interpolated and used in the monitoring of turning points and Stokes lines.\n\n\\begin{figure}[t!]\n \\includegraphics[width=0.49\\textwidth]{NumericalAnalyticContinuation.pdf}\n \\caption{A grid of uniformly spaced points covering a band of the complex plane. The grid points lying over the real axis mark the discretization of physical time. Numerically constructing the Taylor polynomial associated with the frequency function around the point $t_{j}$ allows for the optimal evaluation of $\\Omega_{k}(z_{ij})$ at grid points $z_{ij}$ for which $\\mathrm{Re}\\,z_{ij} = t_{j}$.}\n \\label{Fig. Numerical Analytic Continuation}\n\\end{figure}\n\nWhile the turning points of $\\Omega_{k}(z)$ can be located with the aid of root-finding algorithms designed for multi-valued functions, the problem of determining the Stokes lines sourced by these points requires the numerical integration of an ordinary differential equation. This is evident from the Stokes lines definition ${\\mathrm{Re}\\big[\\Omega_{k}\\mathrm{d}z\\big]=0}$, which implies that, locally, its line element must satisfy ${\\mathrm{d}z \\propto i\/\\Omega_{k}(z)}$. Defining ${t=\\mathrm{Re} \\, z}$ and ${\\tau=\\mathrm{Im} \\, z}$, this condition can be rewritten as\n\\begin{equation}\n\t\\mathrm{d}z = \\mathrm{d}t + i\\,\\mathrm{d}\\tau \\propto \\frac{i}{\\Omega_{k}(z)} \\, .\n\\end{equation}\nTaking the ratio between the matched real and imaginary parts of this proportionality relation leads to the differential equation \n\\begin{equation}\\label{Eq. Stokes ODE}\n\t\\frac{dt}{d\\tau} = \\frac{ \\mathrm{Im} \\, \\Omega_{k}(z) }{ \\mathrm{Re} \\, \\Omega_{k}(z) } \\, \n\\end{equation}\nfor the Stokes line,\nwhich can be numerically integrated from the turning point of interest to yield $t(\\tau)$.\n\nHere is a summary of the minimal set of tasks performed by our algorithm while evolving the physical quantities of interest by one time step:\n\\begin{itemize}\n\t\\item[1.] Take samples of $a(t)$, $H(t)$, and $\\dot{H}(t)$ describing the metric along an interval of the real axis. Over this same interval, sample and interpolate the field energy density $\\rho(t)$ and pressure $P(t)$.\n\t\n\t\\item[2.] Numerically integrate the semi-classical Friedmann equations so as to enlarge the input metric samples $a(t)$, $H(t)$, and $\\dot{H}(t)$ by a time step $\\Delta t$.\n\n\t\\item[3.] For each field mode, generate a sample of the frequency function $\\Omega_{k}(t)$ over the real axis, and numerically extend this function onto the complex plane to obtain $\\Omega_{k}(z)$.\n\t\n\t\\item[4.] Search for complex turning points of each frequency function $\\Omega_{k}(z)$, and numerically trace their corresponding Stokes lines.\n\n\t\\item[5.] If a Stokes line associated with a mode of wavenumber $k$ is found to intersect the real axis, update the real bilinears $\\mathcal{N}_{k}(t)$, $\\mathcal{R}_{k}(t)$ and $\\mathcal{I}_{k}(t)$ accordingly.\n\t\n\t\\item[6.] For each field mode, compute $W_{k}(t)$ and $V_{k}(t)$ up to the optimal truncation order set by the last Stokes line crossing.\n\t\n\t\\item[7.] Gather the results from all previous steps to evolve the input samples for the field energy density $\\rho(t)$ and pressure $P(t)$ by a time step $\\Delta t$.\n\\end{itemize}\n\nIn general, the Stokes lines associated with field modes of comparable wavenumber will cross the real axis within close proximity of one another, giving rise to overlapping particle production events. In order to correctly capture the influence that such events might have on each other, we apply the stepping algorithm outlined above in an iterative fashion. In other words, once the quantities of interest have been forward-stepped up to $t_{j}$, the following iteration backtracks to $t_{0}$ and then proceeds to step the problem up to $t_{j+1} = t_{j} + \\Delta t$ using as sources for the semi-classical Friedmann equations the field energy density and pressure obtained in the previous iteration.\n\nIn summary, our numerical implementation allows for the scale factor and the Stokes geometry to reconfigure themselves with each iteration and thereby construct a self-consistent solution to the back-reaction problem. \n\n\\section{Numerical Results}\\label{Sec. Results}\n\nTo assess the accuracy of our numerical approach, we first neglect back-reaction effects and compare numerical results to known analytic solutions for a quantized scalar field evolving in a closed de Sitter spacetime \\cite{Mottola85}. This case is characterized by a positive cosmological constant $\\Lambda$ and a curvature parameter of $K = 1$, which together lead to a bouncing scale factor evolution\n\\begin{equation}\\label{Eq. de Sitter}\n\ta(t) = \\bar{H}^{-1} \\cosh{\\left(\\bar{H} t\\right)} \\quad \\text{with} \\quad \\bar{H} = \\sqrt{\\Lambda\/3} \\, .\n\\end{equation}\nHere $\\bar{H}$ is the asymptotic value of the Hubble parameter in the infinite future,\n\\begin{equation}\\label{Eq. Hubble de Sitter}\n\t\\lim_{t \\to \\pm \\infty}{H(t)} = \\pm \\bar{H} \\, .\n\\end{equation}\nThis model Universe contracts for $t < 0$, reaches its minimum size at $t = 0$, and subsequently expands for the $t > 0$.\n\nSubstituting Eq.~(\\ref{Eq. de Sitter}) into Eq.~(\\ref{Eq. Mode Frequency}) yields \n\\begin{equation}\\label{Eq. de Sitter Mode Frequency}\n\\Omega_{k}^{2}(t) = \\bar{H}^{2} \\Bigg[ \\bigg(k^{2} - \\frac{1}{4} \\bigg) \\, \\text{sech}^{2} \\! \\left(\\bar{H} t\\right) + \\frac{m^{2}}{\\bar{H}^{2}} + 12\\,\\xi - \\frac{9}{4} \\Bigg] \\, \n\\end{equation}\nfor the mode frequency function.\nAnalytically extending this function to the complex plane, locating its turning points, and tracing its Stokes lines are straightforward. We verify our numerical calculations against these analytic results. For definiteness, we choose a scalar field of mass $m = 0.1\\,M$ which is conformally coupled to the scalar curvature, $\\xi = \\frac{1}{6}$. We set the cosmological constant to $\\Lambda = 3 \\, m^{2}$, so that $\\bar{H} = 1 \\, m$. All dimensional quantities are thus expressed in terms of the field mass.\n\nA comparison between the analytic extension of Eq.~(\\ref{Eq. de Sitter Mode Frequency}) and the numerical analytic continuation produced by our algorithm is displayed in Figure \\ref{Fig. Continuation Error} for the field mode of wavenumber $k = 5 \\, m$. The left panel shows the absolute value of the numerically obtained frequency function, while the right panel exhibits how this result deviates from the analytic expression for $\\Omega_{k}(z)$. In addition to correctly reproducing the function's conjugate pair of zeroes $\\big( z_{0},\\, z_{0}^{\\ast}\\big)$ located in this region, the numerical analytic continuation differs from the analytic value by at most 2\\% in the vicinity of these points. As a result, the Stokes lines which occupy this area of the complex plane can be traced with high fidelity. This is demonstrated in the left panel of Figure \\ref{Fig. Frequency Particle}, where the Stokes lines sourced by the pairs of turning points $\\big( z_{0},\\, z_{0}^{\\ast}\\big)$ and $\\big( z_{1},\\, z_{1}^{\\ast}\\big)$ are superimposed over the numerically obtained frequency function. The effects of each Stokes line on the adiabatic bilinear $\\mathcal{N}_{k}(t)$ are displayed in the right panel of Figure \\ref{Fig. Frequency Particle}, wherein this quantity is tracked as a function of time. Each burst of particle production is prompted by a Stokes line crossing, the first of which occurs as the Universe contracts and the field mode under consideration becomes sub-horizon; while the second burst happens after the bounce, when the mode reverts back to being super-horizon due to the Universe's expansion \\cite{Habib99}. Despite the symmetry between these events, constructive interference expressed by Eq.~(\\ref{Eq. Bilinears Multiple Crossings}) causes more particles to be produced in the second burst. The expected values for the particle number plateaus are indicated by the square markers on the vertical axis, both of which agree well with the numerical curve.\n\n\\begin{figure}[t!]\n \\includegraphics[width=0.49\\textwidth]{ContinuationError.pdf}\n \\caption{A comparison between the numerical analytic continuation of $\\Omega_{k}(z)$ produced by our algorithm and the expected analytic expression for this function in a closed de Sitter spacetime. The field parameters are $m = 0.1 \\, M$, $\\xi = \\frac{1}{6}$, and $k = 5 \\, m$, while the spacetime is characterized by $\\Lambda = 3 \\, m^{2}$ and $K = 1$. The left panel shows the absolute value of the numerically produced frequency function in the vicinity of the pair of conjugate turning points $\\big( z_{0},\\, z_{0}^{\\ast}\\big)$, while the right panel exhibits the relative difference between the analytic and numerical results.}\n \\label{Fig. Continuation Error}\n\\end{figure}\n\n\\begin{figure*}[hbtp]\n \\includegraphics[width=1.0\\textwidth]{FrequencyParticle.pdf}\n \\caption{The numerically traced Stokes geometry associated with the frequency function $\\Omega_{k}(z)$, and the adiabatic particle number evolution $\\mathcal{N}_{k}(t)$ extracted from it. The field parameters are set to $m = 0.1 \\, M$, $\\xi = \\frac{1}{6}$, $k = 5 \\, m$, and $N_{k} = 0$, while the spacetime is characterized by $\\Lambda = 3 \\, m^{2}$ and $K = 1$. The left panel shows the Stokes lines sourced by the pairs of turning points $\\big( z_{0},\\, z_{0}^{\\ast}\\big)$ and $\\big( z_{1},\\, z_{1}^{\\ast}\\big)$ superimposed over the absolute value of the numerically obtained frequency function. The real axis corresponds to the central dashed line. The effects of each Stokes line on the adiabatic particle number $\\mathcal{N}_{k}(t)$ are illustrated on the right panel, wherein this quantity is tracked as a function of time. Each burst of particle production is prompted by a Stokes line crossing, indicated here by the circular markers on the horizontal axis. The expected values for the particle number plateaus featuring in this image are indicated by the square markers on the vertical axis, both of which show very good agreement with the numerically produced curve for $\\mathcal{N}_{k}(t)$. Constructive interference causes more particles to be produced in the second burst.}\n \\label{Fig. Frequency Particle}\n\\end{figure*}\n\\begin{figure*}[hbtp]\n \\includegraphics[width=1.0\\textwidth]{EnergyDensity_NoBackReaction.pdf}\n \\caption{The evolution of every term appearing on the right-hand side of the semi-classical Friedmann Eq.~(\\ref{Eq. First Friedmann Equation}) in a closed de Sitter spacetime evolution, as well as the quantities describing the metric for this spacetime in the absence of back reaction. The field parameters are set to $m = 0.1 \\, M$ and $\\xi = \\frac{1}{6}$, while the spacetime is characterized by $\\Lambda = 3 \\, m^{2}$ and $K = 1$. The bounce starts at $t_{0} = -5 \\, m^{-1}$ with an initial particle distribution given by $\\mathcal{N}_{k}(t_{0}) = 0$. The left panel follows the evolution of $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ (solid line), $H^{2}_{\\mathcal{R},\\,\\mathcal{I}}$ (dot-dashed line), $H^{2}_{\\Lambda}$ (dotted line), and $H^{2}_{K}$ (dashed line). While $H^{2}_{\\mathcal{R},\\,\\mathcal{I}}$ remains negligible throughout, $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ grows exponentially and eventually comes to dominate over all other contributions. The right panels illustrate the scale factor $a(t)$ and Hubble parameter $H(t)$ which describe the de Sitter bounce. Because back-reaction effects are being neglected, the Hubble parameter is just $H^{2} = H^{2}_{\\Lambda} + H^{2}_{K}$.} \n \\label{Fig. Contributions to H2 - No back-reaction}\n\\end{figure*}\n\nBy tracing the Stokes geometry of every field mode, we can also track the evolution of the field energy density as the spacetime evolves. Even though back-reaction effects are being neglected, this quantity shows whether the effects of particle production will eventually become comparable to the contributions from $\\Lambda$ and $K$ which source the background de Sitter spacetime. To that end, we track every term appearing on the right-hand side of the semi-classical Friedmann Eq.~(\\ref{Eq. First Friedmann Equation}), identifying each contribution according to the notation\n\\begin{subequations}\\label{Eq. Contributions to H2}\n\\begin{align*}\n\tH^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}} \\equiv \\frac{\\rho}{3M^2}\\,\\,, \\quad H^{2}_{\\Lambda} \\equiv \\frac{\\Lambda}{3} \\,\\,, \\quad \\text{and} \\quad H^{2}_{K} \\equiv - \\frac{K}{a^{2}} \\, .\n\\end{align*}\n\\end{subequations}\nAdditionally, we define $H^{2}_{\\mathcal{R},\\,\\mathcal{I}}$ as the contribution to the right-hand side of Eq.~(\\ref{Eq. First Friedmann Equation}) which stems solely from terms proportional to the real bilinears $\\mathcal{R}_{k}$ and $\\mathcal{I}_{k}$. The left panel of Figure~\\ref{Fig. Contributions to H2 - No back-reaction} displays the evolution of the above-defined quantities for a bounce that starts at $t_{0} = -5 \\, m^{-1}$ with an initial particle distribution given by $\\mathcal{N}_{k}(t_{0}) = 0$. Being the only true sources in this case, $H^{2}_{\\Lambda}$ and $H^{2}_{K}$ behave in the standard way, acting in concert to produce the de Sitter bounce. Because back-reaction effects are neglected, the Hubble parameter $H$ shown on the right panel of Figure~\\ref{Fig. Contributions to H2 - No back-reaction} is entirely characterized by these two quantities, i.e., $H^{2} = H^{2}_{\\Lambda} + H^{2}_{K}$. On the other hand, the field-related quantities $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ and $H^{2}_{\\mathcal{R},\\,\\mathcal{I}}$ display an interesting behavior which mirrors the result found in Ref.~\\cite{AndersonMottola14}. While $H^{2}_{\\mathcal{R},\\,\\mathcal{I}}$ remains negligible throughout, $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ grows exponentially as the Universe progresses toward the bounce. In other words, the field energy density eventually becomes dominated by $\\mathcal{N}_{k}$ -- the field particle content. Physically, the soaring field energy density is due to the blueshift experienced by particles produced in the contracting phase. As a result, the Universe is filled with relativistic particles which effectively behave as radiation, making the field energy density grow as $\\rho \\propto a^{-4}$. This trend is then reversed in the ensuing expanding phase, during which the field energy density drops rapidly as particles are continuously redshifted.\n\\begin{figure*}[hbtp]\n \\includegraphics[width=1.0\\textwidth]{EnergyDensity_BackReaction.pdf}\n \\caption{The evolution of every term appearing on the right-hand side of the semi-classical Friedmann Eq.~(\\ref{Eq. First Friedmann Equation}) and the quantities describing the metric evolution in a full back-reacting calculation. The field parameters are $m = 0.1 \\, M$ and $\\xi = \\frac{1}{6}$, while the cosmological constant and curvature parameter are $\\Lambda = 3 \\, m^{2}$ and $K = 1$. The closed de Sitter initial conditions are set at $t_{0} = -5 \\, m^{-1}$, along with an initial particle distribution given by $\\mathcal{N}_{k}(t_{0}) = 0$. The left panel follows the evolution of $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ (solid line), $H^{2}_{\\mathcal{R},\\,\\mathcal{I}}$ (dot-dashed line), $H^{2}_{\\Lambda}$ (dotted line), and $H^{2}_{K}$ (dashed line). The exponential growth of $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ effectively fills the Universe with relativistic particles, introducing an instability to the initial de Sitter phase. The right panels illustrate the scale factor $a(t)$ and Hubble parameter $H(t)$ transitioning from a Sitter bounce to a radiation dominated phase. Here the solid lines represent the solutions to the back-reaction problem, while the dashed lines trace the pure de Sitter bounce.} \n \\label{Fig. Contributions to H2 - Back-reaction}\n\\end{figure*}\n\nThe preceding calculations demonstrate that back-reaction effects due to particle production can become dynamically significant in an initially closed de Sitter spacetime. A full account of these effects is shown in Figure~\\ref{Fig. Contributions to H2 - Back-reaction}, using the algorithm for computing back-reaction effects described in the previous Section. In this case, the metric evolution initially matches that of a closed de Sitter spacetime at $t_{0} = -5 \\, m^{-1}$, while the initial particle distribution is given by $\\mathcal{N}_{k}(t_{0}) = 0$. These initial conditions self-consistently satisfy the semi-classical Friedmann equations at the initial time $t_{0}$ within our approximations. As in the case without back reaction, the quantity $H^{2}_{\\mathcal{R},\\,\\mathcal{I}}$ remains sub-dominant throughout the evolution, while $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ grows exponentially as newly-created particles are continuously blueshifted. Since they quickly become relativistic, these particles behave as an additional radiation-like component, destabilizing the initial de Sitter phase. This is illustrated in the right panels of Figure~\\ref{Fig. Contributions to H2 - Back-reaction}, where the scale factor and Hubble parameter can be seen transitioning from a de Sitter bounce to a radiation-dominated behavior. The contributions from the regularized vacuum terms discarded in our approximations remain negligible at all times. We stop the numerical integration at $t = -1.3 \\, m^{-1}$, since beyond this time the Hubble parameter becomes of order $H \\simeq M^{-1}$, invalidating the semi-classical picture of gravity on which our calculations rely.\n\nThe de Sitter bounce is not always disrupted by particle production. For sufficiently low values of the field mass, the bounce is merely delayed. Figure~\\ref{Fig. Contributions to H2 - Back-reaction 2} illustrates a near-limiting case with ${m = 0.0145 \\, M}$ for which an initial de Sitter evolution is still driven toward a radiation dominated phase. The contributions due to particle production $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ only come to dominate over the combined $H^{2}_{\\Lambda}$ and $H^{2}_{K}$ near the bounce at $t = 0 \\, m^{-1}$. For field masses $m \\lesssim 0.0142 \\, M$, the negative curvature contributions $H^{2}_{K}$ neutralize the growth of $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ for long enough to preserve the bounce. The resulting bounce is pushed to a slightly later time and occurs at a smaller value of the scale factor.\n\n\\begin{figure*}[hbtp]\n \\includegraphics[width=1.0\\textwidth]{EnergyDensity_BackReaction_2.pdf}\n \\caption{The evolution of every source term featuring on the right-hand side of the semi-classical Friedmann Eq.~(\\ref{Eq. First Friedmann Equation}) and the quantities describing the metric evolution in a full back-reacting calculation. The field parameters are set to $m = 0.0145 \\, M$ and $\\xi = \\frac{1}{6}$, while the cosmological constant and curvature parameter are characterized by $\\Lambda = 3 \\, m^{2}$ and $K = 1$. The closed de Sitter initial conditions are set at $t_{0} = -5 \\, m^{-1}$, along with an initial particle distribution given by $\\mathcal{N}_{k}(t_{0}) = 0$. The left panel follows the evolution of $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ (solid line), $H^{2}_{\\mathcal{R},\\,\\mathcal{I}}$ (dot-dashed line), $H^{2}_{\\Lambda}$ (dotted line), and $H^{2}_{K}$ (dashed line). The growth of $H^{2}_{\\mathcal{N},\\,\\mathcal{R},\\,\\mathcal{I}}$ fills the Universe with just enough relativistic particles to destabilize the initial de Sitter phase. The right panels illustrate the scale factor $a(t)$ and Hubble parameter $H(t)$ transitioning from a de Sitter bounce to a radiation dominated phase. Here the solid lines represent the solutions to the back-reaction problem, while the dashed lines trace the pure de Sitter bounce. Had the field mass been set to a value $m \\lesssim 0.0142 \\, M$, a bounce would still take place, albeit at a slightly later time and for a smaller value of the scale factor.}\n \\label{Fig. Contributions to H2 - Back-reaction 2}\n\\end{figure*}\n\n\\section{Discussion}\\label{Sec. Discussion}\n\nThe back-reaction problem addressed in this work imposes several technical hurdles which have resisted a satisfactory solution for decades. These difficulties stem primarily from the necessity to control the divergent nature of the vacuum energy. Adiabatic regularization accomplishes this at the cost of increasing the problem's complexity. As a result, ambiguities arise in the specification of initial conditions and in the value of physical quantities when particle production is rapid, and computationally the problem becomes susceptible to potential numerical instabilities. In this work we have shown that these issues can be circumvented in scenarios dominated by particle production. Our approach relies on a particular choice of adiabatic mode functions which isolate the vacuum contributions into a separate covariantly conserved component of the total stress-energy. In regimes dominated by particle production, this vacuum component is sub-dominant and can be discarded in its entirety. By definition, the remaining covariantly conserved portion of the stress-energy dominates, as it encapsulates the effects of particle production. This component can be expressed in terms of the particle number density as described by Berry's universal form, resolving the ambiguity in physical quantities, and computed from the analytic continuation of each mode's frequency function onto the complex plane (Figures \\ref{Fig. Frequency Particle} and \\ref{Fig. Contributions to H2 - No back-reaction}). The resulting stress-energy is a calculable source term for the semi-classical Friedmann equations, and can be used to obtain a numerical solution to the back-reaction problem. We have performed this calculation for an initially closed de Sitter spacetime, demonstrating that the effects of particle production in this scenario can become strong enough to drive the cosmic evolution into a radiation-dominated phase (Figures \\ref{Fig. Contributions to H2 - Back-reaction} and \\ref{Fig. Contributions to H2 - Back-reaction 2}). Our results illustrate the reliability of our numerical implementation, and open the possibility of the systematic investigation of cosmological scenarios dominated by quantum particle production.\n\nOn a technical level, our method relies on some previous knowledge of the Stokes geometry associated with the spacetime evolution. For the case studied in this work, all Stokes lines are sufficiently separated from each other so that Berry's universal form for particle production applies without corrections. In general, however, the spacetime evolution might result in near-lying Stokes lines for which higher-order Stokes corrections are required for an accurate description of particle production. Although not included in this work, such corrections are well-documented in the literature \\cite{BerryHowls94, HowlsLongmanDaalhuis04, HowlsDaalhuis12} and could in principle be added to our numerical implementation. More fundamentally, our method is based on a well-defined semi-classical notion of particle. Mathematically, this notion is tied to the existence of a phase-integral expansion for the field mode functions. Such a representation can always be constructed as long as $\\left| \\varepsilon_{k,\\,0} \\right| \\ll 1$. Physically, this requirement typically translates to an approximate bound on the Hubble rate $H \\lesssim m$ set by the mass of the field under consideration. Nonetheless, some scenarios exist for which $\\left| \\varepsilon_{k,\\,0} \\right| \\ll 1$ is satisfied even when $H > m$.\n\nQuantum backreaction is potentially important in models of the very early Universe. Quantum particle production is actually quite familiar in the context of inflation, as it provides the standard mechanism for the generation of perturbations in an inflating spacetime \\cite{Starobinsky79, Allen88, Sahni90, Mukhanov92, Souradeep92, Glenz09, Agullo11}. Interestingly, it has been suggested that these same ideas could be applied to the problem of driving inflation itself \\cite{Prigogine89, Calvao92, Lima96, Abramo96, Gunzig98, Lima14, HaroPan16}. Indeed, a phase of accelerated expansion can result if particles are produced at a high enough rate. A time derivative of the usual Friedmann Equation $H^{2}(t) = \\frac{1}{3}M^{-2} \\rho(t)$ shows that an accelerating expansion ${\\ddot a} > 0$ occurs when\n\\begin{equation}\\label{Eq. Accelerated Condition}\n\t\\dot{\\rho}(t) > -\\frac{2}{\\sqrt{3}} \\, M^{-1} \\rho^{3\/2}(t) \\, .\n\\end{equation}\nSuch a scenario has the potential to sidestep some of the conceptual problems of the standard inflationary paradigm. For instance, it has been argued that standard inflation cannot generically start in a patch which is smaller than the cosmological horizon without violating either causality or the weak energy condition \\cite{VachaspatiTrodden99, BereraGordon01}. However, if inflation is initially driven by an increasing energy density due to particle production, the weak energy condition {\\it is} effectively violated. Therefore, inflation driven by such a mechanism could generically start in small patches contained within the cosmological horizon without violating causality. Inflation driven by particle production would also clarify the meaning of the inflaton effective potential by making manifest the high mass-scale physics it represents.\n\nThe same conditions which lead to quantum particle production can also result in particle annihilation. If sufficiently pronounced, this effect can drive a contracting spacetime toward a bounce phase. Indeed, it follows from the cosmological continuity equation that $\\rho(t) + P(t) < 0$ provided the particle annihilation rates are high enough to cause the field energy density to decrease as the Universe contracts. In other words, the null energy condition is effectively violated, making $\\dot{H}(t) > 0$ according to the Friedmann equations \\cite{IjjasSteinhardt18, Ijjas16}. Thus, a classical bounce can emerge provided enough energy density is sequestered by quantum particle annihilation during a phase of cosmological contraction. If realized, such a mechanism could provide a natural description for cosmological bounce scenarios which does not require new physics. Also, successful bounces require constraints on high mass-scale quantum fields, so that quantum back-reaction does not push the contracting phase into a radiation crunch, as with the example solved in this paper.\n\nAnother possibly interesting effect is the production of a relativistic condensate in the early universe. Under certain circumstances, quantum particle production can lead to large occupation numbers for some scalar field modes, representing condensate formation. This phenomenon could lead to additional interesting phenomenology \\cite{AragaoRosa80, ParkerZhang91, ParkerZhang93}. \n\nA number of technical questions remain to be answered. Fermion fields require more complex calculations than scalar fields, and may present some different physics \\cite{Landete13, Landete14}. How to handle interacting fields remains an open question, and multiple fields offer additional possibilities \\cite{Ringwald87, CooperMottola87, PazMazzitelli88, Habib96, Cooper97, MolinaParis00}. We have a long road to travel before the range of interesting early-Universe dynamical scenarios driven by quantum particle production has been fully explored. \n\n\n\\begin{acknowledgments}\nF.Z. thanks D. Boyanovsky, S. Habib, P. Anderson, and E. Mottola for helpful discussions. F.Z. acknowledges support from the Andrew Mellon Predoctoral Fellowship and the A\\&S PITT PACC Fellowship. The authors have been partly supported by the National Science Foundation under the grant AST-1312380. This work made use of many community-developed or community-maintained software packages, including (in alphabetical order): Matplotlib \\cite{Matplotlib}, NumPy \\cite{NumPy}, and SciPy \\cite{SciPy}. Bibliographic information was obtained from the NASA Astrophysical Data System.\n\\end{acknowledgments}\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzcpgl b/data_all_eng_slimpj/shuffled/split2/finalzzcpgl new file mode 100644 index 0000000000000000000000000000000000000000..d6a39eb749cc08ef02a3ff57efe030ffc88bc742 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzcpgl @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\n\n\nConsider a set of combinatorial objects $\\{\\alpha\\}$ with statistics $wt(\\alpha)$ and $t(\\alpha)$, thinking of $wt(\\alpha)$ as the primary descriptor. Let $G(z,q)$ be its two-variable generating function, that is, if $p(n,k)$ is the number of objects $\\alpha$ with $wt(\\alpha) = n$ and $t(\\alpha) = k$, \n \\[\n G(z,q) = \\sum_{n,k \\in \\mathbb{N} \\bigcup \\{0\\}} p(n,k) q^n z^k\\, \\text{.}\n \\]\n \nAn important example occurs when the generating function\n $G(z,q)$ has the Euler product form \n\\begin{equation}\\label{eq:product}\nG(z,q) = \\prod_{i \\geq 1} \\frac{1}{(1-z q^i)^{a_i}}\n\\end{equation}\nwith $a_i \\in \\mathbb{Z}^+ \\bigcup \\{0\\}$. \nWe let $F_n(z)$ denote the $q^n$ coefficient of $G(z,q)$ which is a polynomial in $z$, so\n\\[\nG(z,q) = \\sum_n F_n(z) q^n, \\quad \\, F_n(z) = \\sum_k p(n,k) z^k \\, .\n\\]\n\nFor $a_i = 1$ this is the generating function for partitions of weight $wt(\\alpha) = n$ with number of parts $t(\\alpha) = k$. It is a well-known fact in partition theory that $p(n,k)$ for $k$ nearly equal to $n$ has a value independent of $n$: $p(n,n-b)$ is the number of partitions of $b$ for $b\\leq\\frac{n}{2}$. A similar result holds for plane partitions of $n$ indexed by their trace \\cite[p. 199]{Andrews_book} or \\cite[Corollary 5.3]{Stanley}.\n\nThis phenomenon, which we call {\\sl stabilization,} is widespread in generating functions of combinatorial interest, even those of greater complexity. The purpose of this paper is to describe this behavior in more general cases, and consider some illustrative examples and variations. \nWe found the polynomial framework to be well suited to these problems rather than a direct approach. The arguments should be adaptable to a wide variety of cases.\n\nWe would like to thank the anonymous referee for a careful readthrough: noting typos, improving exposition, and suggesting occasional strengthening of theorems for which we had been hesitant to extend our reach. This article is substantially improved from their efforts.\n\n\\section{ Basic Infinite Product Generating Functions }\n\nLet $G(z,q)$ be an infinite product generating function of the form\n\\begin{equation}\\label{eq:new_genfct}\nG(z,q) = \\sum_{n=1}^\\infty F_n(z) q^n = \\prod_{j=1}^{\\infty} \\frac{1}{(1-z^{b(j)} q^{c(j)})^{a_j}}.\n \\end{equation}\n where we assume that the number of $j$ for which $c(j) = t$ for any $t$ is finite, so that the series converges. We find that if the $c$ grow sufficiently faster than the $b$, the upper ends of the $F_n(z)$ stabilize, to the coefficients of a single-variable generating function which we can give.\n \n Let ${\\mathcal F}$ denote the set of all nonnegative integer sequences with finite support.\n For ${\\bf e} = (e_1,e_2,\\dots) \\in {\\mathcal F}$\nset\n \\[\n \\mu(\\mathbf{e}) = (c(1)^{e_1}c(2)^{e_2}\\cdots),\\quad\n \\nu(\\mathbf{e}) = (b(1)^{e_1}b(2)^{e_2}\\cdots)\n \\]\nto denote the partitions with parts $c(j)$ (resp. $b(j)$) appearing $e_j$ times. \nA direct expansion of the generating functions yields\nan explicit form for the polynomial $F_n(z)$:\n\\begin{eqnarray}\\label{eq:take_away}\nF_n (z)\n &= &\n\\sum_{k=0}^n z^k \\, \\sum_{{\\mu(\\mathbf{e}) \\vdash n} \\atop {\\nu(\\mathbf{e}) \\vdash k}} \\prod_{i \\geq 1} \n \\binom{ a_i+e_i-1}{e_i}\n \\end{eqnarray}\n We will compare the coefficients of $F_n(z)$ with those of the expansion of\n \\[\n \\prod_{i=1}^\\infty\n \\frac{1}{ (1-z^{b(i)})^{a_i}} =\n \\sum_{k=0}^\\infty z^k \\, \\sum_{ \\nu( {\\bf e}) \\vdash k} \\, \\prod_{i\\geq 1} \\binom{ a_i + e_i -1}{ e_i} .\n \\]\n \n\n\\begin{theorem}\n\\label{thm:stabliztion2}\n Suppose that the exponents satisfy $a_1 = 1 = b(1) = c(1)$.\n If there exists a positive integer $m \\geq 2$ such that \n \\begin{equation}\\label{eq:stable2}\n m \\cdot b(j) \\leq c (j), \\quad j \\geq 2 ,\n \\end{equation}\n then \n \\\\\n (a)\n for $k > n\/m $, $\\left[z^k \\right] F_n(z) = \\left[z^{k+1} \\right] F_{n+1} (z)$.\n \\\\\n (b) \n If $c(j) - b(j) > 0$ for all $j>1$ and the set of $j$ for which $c(j) - b(j)$ takes a given value is finite for any fixed difference, \n then for $\\ell \\leq \\lfloor n\/m \\rfloor$, \n \\[\n \\left[ z^{n-\\ell}\\right]F_n(z) = \\left[z^{\\ell}\\right] \\prod_{j \\geq 2} \\frac{1}{(1-z^{c(j)-b(j)})^{a_j}} .\n \\]\n\\end{theorem}\n\\begin{proof}\n(a)\nFrom the explicit form (\\ref{eq:take_away}) of the polynomial $F_n(z)$, we can expand it as\n\\begin{eqnarray} \nF_n (z)\n &= &\n\\sum_{k=0}^n z^k \\, \\sum_{{\\mu(\\mathbf{e}) \\vdash n} \\atop {\\nu(\\mathbf{e}) \\vdash k}} \\prod_{i \\geq 1} \n \\binom{ a_i+e_i-1}{e_i}\n \\nonumber\n\\\\ \n&=& \\sum_{k=0}^n \\sum_{e_1=0}^k z^{b(1) e_1}\n \\sum_{{\\mu^-(\\mathbf{e}) \\vdash n - c(1)e_1} \\atop {\\nu^-(\\mathbf{e}) \\vdash k - b(1)e_1}} z^{k-b(1)e_1}\n \\prod_{i \\geq 2} \\left( {{a_i+e_i-1} \\atop {e_i}} \\right) .\n\\nonumber\n\\end{eqnarray}\n\n\nNow if the integer sequence ${\\bf e}$ gives a contribution to $[z^{k+1}]F_{n+1}(z)$ and $e_1>0$, then we define \n${\\bf e}'$ \nas the integer sequence \nall of whose terms agree with ${\\bf e}$ except for $j=1$ where we set $e_1'=e_1-1$.\nIn this way, we obtain all the possible terms contributing to $[z^k]F_n(z)$.\nConversely, any term for $[z^k]F_n(z)$ gives a contribution to $[z^{k+1}]F_{n+1}(z)$ by simply adding\n1 to its first component.\n\nThe result reduces to showing that any contribution to $[z^{k+1}]F_{n+1}(z)$ indexed by ${\\bf e}$\nmust have $e_1>0$.\n\nWe introduce the notation for the modified partitions \n\\[\n \\mu^{-}(\\mathbf{e}) = (c(2)^{e_2} c(3)^{e_3} \\cdots),\\quad\n \\nu^{-}(\\mathbf{e}) = (b(2)^{e_2} b(3)^{e_3} \\cdots) .\n \\]\nNow assume that the partition with $\\mu({\\bf e}) \\vdash n$ and $\\nu({\\bf e}) \\vdash k$ gives a contribution to $[z^{k+1}]F_{n+1}(z)$ but\n $e_1 = 0$. So $\\mu^-(\\mathbf{e}) \\vdash n$ and $\\nu^-(\\mathbf{e}) \\vdash k$, but if $c(j) \\geq m \\cdot b(j)$ for all $j \\geq 2$, \n then $\\vert \\mu^-({\\bf e}) \\vert \\geq m \\vert \\nu^-({\\bf e}) \\vert$.\nHence if $ k > \\frac{n}{m}$, we find that $\\vert \\mu^{-}({\\bf e}) \\vert > n$, a contradiction. Thus all terms in both expansions are the same, and so the coefficients are equal. This proves part (a).\n\nFor part (b), we begin by forming a new partition as follows.\nFirst subtract $b(j)$ from each $c(j)$ and consider the partition\n$\\lambda(\\mathbf{e}) = ((c(2)-b(2))^{e_2}\\dots)$. This removes exactly the amount $\\vert \\nu^{-}(\\mathbf{e})\\vert$ from $\\vert \\mu^{-}(\\mathbf{e})\\vert$, so\n\\begin{eqnarray}\nF_n (z)\n&=& \\sum_{k=0}^n \\sum_{e_1=0}^k z^{k}\n \\sum_{{\\lambda(\\mathbf{e}) \\vdash n - k} \\atop {\\nu^-(\\mathbf{e}) \\vdash k - b(1)e_1}} \\prod_{i \\geq 2} \\left( {{a_i+e_i-1} \\atop {e_i}} \\right)\n\\nonumber\n\\\\\n&=& \\sum_{k=0}^n z^{k}\n \\sum_{{\\lambda(\\mathbf{e}) \\vdash n - k} \\atop {\\vert \\nu^-(\\mathbf{e}) \\vert \\leq k}} \\prod_{i \\geq 2} \\left( {{a_i+e_i-1} \\atop {e_i}} \\right) . \\nonumber\n\\end{eqnarray}\n\nBy hypothesis, the parts of $\\lambda$ satisfy $c(j) - b(j) \\geq (m-1)b(j)$. \nIf $k = n-\\ell > \\lfloor n\/m \\rfloor$, then $\\ell < \\lfloor n \\frac{m-1}{m} \\rfloor$. \n If $\\lambda(\\mathbf{e}) \\vdash n-k = \\ell$, then $\\nu^-(\\mathbf{e}) \\leq \\frac{\\ell}{m-1} \\leq k$. \n Thus the sum runs over all $\\mathbf{e}$ for which $\\lambda(\\mathbf{e}) \\vdash \\ell$ with parts $c(j) - b(j)$. But this is exactly the coefficient of $z^{\\ell}$ in the expansion claimed: \n \\[\n \\left[ z^{n-\\ell}\\right]F_n(z) = \\left[z^{\\ell}\\right] \\prod_{j \\geq 2} \\frac{1}{(1-z^{c(j)-b(j)})^{a_j}} \\, \\text{ .}\n \\]\n\\end{proof}\n\nBy the technique of proof, we have a slight improvement in a special case of hand enumerators of prefabs\n (\\cite{Bender_Goodman}, \\cite[page 92]{Wilf}).\n\n\\begin{corollary}\nSuppose $a_1=1$.\nIf $F_n(z) = [q^n] \\prod_{j=1}^\\infty (1-zq^j)^{- a_j}$,\n then\n for $k \\geq n\/2$, $\\left[z^k \\right] F_n(z) = \\left[z^{k+1} \\right] F_{n+1} (z)$.\n\\end{corollary}\n\nBy the proof, we find that the generating function for the stabilized coefficients gives an upper bound outside the range of stability.\n\n\\begin{corollary}\nWith the hypotheses of the Theorem,\n\\[\n \\left[ z^{n-\\ell}\\right]F_n(z) \\leq \\left[z^{\\ell}\\right] \\prod_{j \\geq 2} \\frac{1}{(1-z^{c(j)-b(j)})^{a_j}} ,\n \\quad 0 \\leq \\ell \\leq n. \n \\]\n\\end{corollary}\n\nVariants of stabilization exist in several guises. If the $b$ grow faster than the bound of the previous theorem, we find that the smaller end of the polynomials stabilize instead of the larger (and $b$ and $c$ staying within a given ratio range will permit both phenomena). We also note that it is possible for the larger coefficients of a sequence of polynomials to stabilize in periods, i.e., the coefficients match those of a polynomial every 2 or more steps further along.\n\n\\begin{theorem}\\label{thm:stabilization2}\n(a)\nLet $\\{ b(j)\\}$ and $\\{ c(j)\\}$ be two strictly increasing sequences of integers, positive except that \n$b(1)=0$. \n Let $F_n(z) = [q^n]\\prod_{j=1}^\\infty (1-z^{ b(j) } q^{ c(j) })^{ - a_j }$ with $a_1=1$.\nIf there exists a positive integer $m $ such that for all $j \\geq 2$\n \\[\n (m+1) b(j) > c(j) ,\n \\]\nthen \n\\[\n[z^k] F_n(z) = \\, [ z^k] F_{n+ c(1)}(z), \\quad 0 \\leq k \\leq n\/(m+1).\n\\]\n(b)\nLet $a_1=0, a_2=1$. \nIf\n$F_n(z) = [q^n]\\prod_{j=2}^\\infty (1-z q^{ j})^{ - a_j }$, then\n$\\deg(F_n) = \\lfloor n\/2 \\rfloor$\nand\n\\[\n[z^k] F_n(z) = \\, [z^{k+1}] F_{n+2}(z) , \\quad k \\geq n\/3 .\n\\]\n\\end{theorem}\n\\begin{proof}\nLet ${\\bf e}$ be a sequence of non-negative integers with finite support.\nBy (\\ref{eq:take_away}), for the coefficients $[z^k]F_n(z)$ and $[z^k] F_{n+c(1)}(z)$, we need to consider the following two sets of \nfinitely supported nonnegative integer sequences:\n\\begin{eqnarray*}\nS_n &=&\n \\{ {\\bf e} \\in {\\mathcal F} : | \\mu({\\bf e}) | = n, \\, | \\nu({\\bf e} ) | =k \\}, \n \\\\\nS_{n+c(1)} &=&\n \\{ {\\bf f} \\in {\\mathcal F} : | \\mu({\\bf f}) | = n+c(1), \\, | \\nu({\\bf f} ) | =k \\} .\n\\end{eqnarray*}\nWe construct a bijection between these two sets. Given ${\\bf e} \\in S_n$, we take the\ncorresponding ${\\bf f}$ with $f_i=e_i$ for $2 \\leq i$ since $b(1)=0$.\nNext write out $ | \\mu({\\bf e}) | = n$\n and $ | \\mu({\\bf f}) | = n+c(1)$:\n\\[\nn= \\sum_{j \\geq 1} c(j) e_j, \\quad n+c(1) = c(1) f_1 + \\sum_{j \\geq 2} c(j) e_j .\n\\]\nThe term $f_1$ uniquely determines a preimage $e_1$ \nprovided $f_1>0$. But $f_1$ must be positive for $k \\leq n\/(m+1)$ since\n\\[\n\\sum_{j\\geq 2} c(j) f_j = \\sum_{j \\geq 2} c(j) e_j < (m+1) \\sum_{j \\geq 2} b(j) e_j = (m+1)k \\leq n.\n\\]\nFinally, the coefficients themselves agree; that is, $[z^k] F_n(z) = [z^k] F_{n+c(1)}(z)$:\n\\[\n \\prod \\left\\{ \\binom{ a_i + e_i -1}{ e_i} : {\\bf e} \\in S_n \\right\\} =\n \\prod \\left\\{ \\binom{ a_i + f_i -1}{ f_i} : {\\bf f} \\in S_{n+c(1)} \\right\\} \n\\]\nsince the above binomial coefficients are all equal for $i \\geq 2$ and when $i=1$ they both reduce to $1$ since $a_1=1$.\n\n\nFor part (b), let\n$S_n = \\{ {\\bf e } \\in {\\mathcal F} :\n\\mu({\\bf e}) \\vdash n, \\nu({\\bf e} ) \\vdash k\\}$\nwhile\n$S_{n+2}=\n \\{ {\\bf f } \\in {\\mathcal F} :\n\\mu({\\bf f}) \\vdash n+2, \\nu({\\bf f} ) \\vdash k+1\\}$.\nWe can construct a bijection between these two sets as in part (a) provided $f_2 >0$ when $k \\geq n\/3$.\nAssume that $f_2=0$ is possible. Then $\\sum f_j = k+1$ while $\\sum_{j \\geq 3} j f_j = n+2$. On the other hand,\n$3(k+1) \\geq \\sum_{j \\geq 3} j f_j $ which yields a contradiction.\n\\end{proof}\n\n\n\n \n \\section{Partitions With Prescribed Subsums}\\label{section:subsums}\n \n \n \nFix a positive integer $m$ and integer $i$ so $1 \\leq i \\leq m$. Canfield-Savage-Wilf \\cite[Section 3]{CSW} introduced the generating function\n\\[\nG_{m,i}(z,q) = \\prod_{j=1}^\\infty \\, \n\\prod_{b=1}^{i-1} \\frac{1}{ 1- z^{j-1} q^{ (j-1)m+b}} \\, \\prod_{b=i}^m \\frac{1}{ 1- z^{j} q^{ (j-1)m+b}}\n\\]\nto describe partitions with prescribed subsums. They let\n$\\Lambda_{m,i}(n,k)$ be the number of partitions $\\lambda=(\\lambda_1,\\cdots, \\lambda_n)$ of $n$\nsuch that the sum of those parts $\\lambda_j$ whose indices $j$ are congruent to $i$ modulo $m$ is $k$;\nthat is, \n\\[\n\\sum_{ j : j \\equiv i \\, ({\\rm mod} \\, m)} \\lambda_j = k .\n\\]\nThen they found\n\\[\n\\sum_{n,k \\geq 0} \\Lambda_{m,i}(n,k) z^k q^n = G_{m,i}(z,q).\n\\]\n\nWe begin by recovering a result for \n$\\Lambda_{2,2}(n,k)$ in \\cite[Theorem 1]{CSW} and \\cite{Yuri}\nby reformulating it in terms of the generating function\n$G_{2,2}(z,q)$ and stabilization of polynomial coefficients.\n\n\\begin{proposition}\\label{prop:n-mk}\nLet $m \\geq 2$ and $1 \\leq b < m$. Let $G(z,q)= \\prod_{j=1}^\\infty (1-z^{j-1} q^{(j-1)m+b} )^{-1}$\nand $A_n(z) = [q^n] G(z,q)$. Then for $ 0 \\leq k \\leq n\/(m+1)$ we have\n\\begin{enumerate}\n\\item\\quad\n$[z^k] A_n(z) = [z^k] A_{n+b}(z)$,\n\\item\\quad\nif $n-mk$ is not divisible by $b$, then\n$[z^k] A_n(z) = 0$,\n\\item\\quad\nif $n-mk$ is divisible by $b$ and $bk \\leq n\/(m+1)$, then\n$[z^k] A_n(z) = p(k)$.\n\\end{enumerate}\n\\end{proposition}\n\\begin{proof}\nThe first part is a direct consequence of the last theorem.\n\nFor part (2), in Theorem \\ref{thm:stabilization2} let $c(j) = m(j-1)+b$ and $b(j)=j-1$. When all $a_j =1$, $[z^j]F_n(z)$ is the number of all ${\\bf e} \\in {\\mathcal F}$ such that $\\mu({\\bf e}) \\vdash n$ and $\\nu({\\bf e}) \\vdash k$.\nConsider\n\\begin{eqnarray}\\label{eq:n-mk}\nn\n&=&| \\mu({\\bf e})| = \\sum_{j \\geq 1} e_j [ m(j-1) + b] \n\\nonumber\n\\\\\n&=&\nm \\sum_{j \\geq 1} e_j (j-1) + b \\sum_{j \\geq 1} e_j \n\\nonumber\n\\\\\n&=&\nmk + b \\sum_{j \\geq 1} e_j .\n\\end{eqnarray}\nHence $n-mk$ must be divisible by $b$ for any nonzero choice of {\\bf e}.\n\n\nFor part (3),\nassume $0 \\leq bk \\leq n\/(m+1)$ and that $n-mk$ is divisible by $b$. \nLet {\\bf e} be any solution to $\\nu( {\\bf e})=k$. Note that this does not give\na constraint for the choice of $e_1$. On the other hand, the choice of $e_1$ by (\\ref{eq:n-mk}) must be\n\\[\ne_1 = (n-mk)\/b - \\sum_{j\\geq 2} e_j, \\quad e_1 \\geq 0\n\\]\nto yield $\\mu ({\\bf e}) = n$. \nHence, $0\\leq [z^k] F_n(z) \\leq p(k)$.\nFinally,\nthe inequality $0 \\leq bk \\leq n\/(m+1)$ shows that $e_1\\geq 0$ always holds. Hence, \n$[z^k] F_n(z)=p(k)$.\n\\end{proof}\n\n\n\n \\begin{theorem}\n Let $G_A(z,q)=\\prod_{k=1}^\\infty ( 1-z^{k} q^{2k-1} )^{-1}$\n and $G_B(z,q)= \\prod_{k=1}^\\infty ( 1-z^{k} q^{2k} )^{-1}$, so $G_{2,2}(z,q) = G_A(z,q) G_B(z,q)$. Then the coefficients of the polynomials\n $F_n(z)= [q^n] G_{2,2}(z,q)$ satisfy\n \\[\n [z^{k}] F_n(z) = [z^{k}] F_{n+1}(z) = \\sum_{\\ell=0}^k p(\\ell) p(k-\\ell), \\quad 0 \\leq k \\leq n\/3,\n \\]\n where $p(\\ell)$ is the number of partitions of $\\ell$, as usual.\n \\end{theorem}\n \\begin{proof} \n Let $A_n(z)=[q^n]G_A(z,q)$ and $B_n(z) = [q^n]G_B(z,q)$. The generating function $G_B(z,q)$ has the explicit expansion\n \\[\n G_B(z,q) = \\sum_{j=0}^\\infty p(j) z^j q^{2j} \n \\]\n so $B_{2k+1}(z)=0$ while $B_{2k}(z) = p(k) z^k$.\n We also know by Proposition \\ref{prop:n-mk} that\n \\[\n [z^j] A_n(z) = p(j), \\quad 0 \\leq j \\leq n\/3.\n \\]\n Next we have that\n \\[\n F_n(z) = \\sum_{\\ell=0}^{n\/2} A_{n-2\\ell}(z) B_{2\\ell}(z) = \\sum_{\\ell=0}^{n\/2} p(\\ell) z^\\ell A_{n-2\\ell}(z)\n \\]\n Examine the coefficient $[z^k]F_n(z)$ for $0 \\leq k \\leq n\/3$:\n \\[\n [z^k]F_n(z) = [z^k] \\sum_{\\ell=0}^{n\/2} p(\\ell) z^\\ell A_{n-2\\ell}(z)\n =\\sum_{\\ell=0}^{n\/2} p(\\ell) \\, [z^{k-\\ell} ] A_{n-2\\ell}(z)\n \\]\n Since $k-\\ell \\leq (n-2\\ell)\/3$ for $0\\leq \\ell \\leq k$ and $k\\leq n\/3$,\n we find $[z^{k - \\ell}] A_{n - 2 \\ell}(z) = p(k - \\ell)$.\n We conclude that\n \\[\n [z^k]F_n(z)= \\sum_{\\ell=0}^{k} p(\\ell) \\, p( k - \\ell) .\n \\]\n \\end{proof}\n \n\n\n \n\nIn order to investigate more general cases for the\n polynomials $F_n(z) = [q^n] G_{m,i}(z,q)$ it is convenient to rewrite the generating function as\n\\begin{equation}\\label{eq:new_eq}\nG_{m,i}(z,q)=\n\\left( \\prod_{b=1}^{i-1} \\frac{1}{1-q^b} \\, \\right) \\,\n\\prod_{a=1}^\\infty \\, \\prod_{d=0}^{m-1} \\frac{1}{ 1-z^a q^{ i+ (a-1) m + d}}.\n\\end{equation}\nWe wish to find a useful form for the coefficient $[z^j] F_n(z)$. As usual, we have\n\\begin{eqnarray*}\n[q^n] G_{m,i}(z,q)\n&& = \n\\sum_{s=0}^n [q^{n-s}] \\left( \\prod_{b=1}^{i-1} \\frac{1}{1-q^b} \\, \\right) \\, \n[q^s] \\prod_{a=1}^{\\infty} \\prod_{d=0}^{m-1} \\frac{1}{ 1- z^{a} q^{i+ (a-1)m+d}}\n\\\\\n&&= \\quad\n\\sum_{s=0}^n [q^{n-s}] \\left( \\prod_{b=1}^{i-1}\\, \n \\sum_{ f_b=0}^\\infty \\, \\left(q^b \\right)^{f_b} \\, \\right) \n\\\\\n&&\n\\qquad\n\\times \\quad\n[q^s] \\prod_{a=1}^{\\infty} \\prod_{d=0}^{m-1} \\, \\sum_{ f_{(a,d)}=0}^\\infty \n\\left(z^{a} q^{ i+(a-1)m+d} \\right)^{f_{(a,d)}} \n\\end{eqnarray*}\nA typical term of the above expansion is indexed by a pair of partitions $(\\rho, \\mu)$\nwhere $\\rho \\vdash n-s$ where each part of $\\rho$ is $ s \/(i+1)$, then $\\mu$ must have at least one part of size $i$.\n\\end{lemma}\n\\begin{proof}\nWe now assume that $m \\geq i+2$ and that the part $i$ does not appear in $\\mu$.\nIn particular, for any part of $\\mu$ of the form $i +d$, with $a=1$, $d$ must be strictly positive.\nThen we have a refinement of the above bounds:\n\\begin{eqnarray*}\ns &\\geq&\ni r + m (j-r) + \\sum_{a,d} d e_{\\mu(a,d)}\n\\geq\ni r + (i+2) (j-r) + \\sum_{ d} d e_{1,d}\n\\\\\n&\\geq&\ni r + (i+2) (j - r) + (r-t)\n\\\\\n&\\geq&\n i r + (i+1) (j - r) + t + (r-t)\n \\\\\n &=&\n (i+1) [ r + (j-r)] = (i+1) j.\n\\end{eqnarray*}\nHence the partition $\\mu$ must have $i$ as a part; otherwise, $(i+1)j >s$ which contradicts our assumption.\n\\end{proof}\n\n\nWith these preparations, we will show that\n\n\\begin{theorem}\\label{thm:subsums}\nLet $m \\geq 1$ and $1 \\leq i \\leq m$. Set $F_n(z) = [ q^n] G_{m,i}(z,q)$ where\n$G_{m,i}(z,q)$ is given by (\\ref{eq:new_eq}).\n If $m>i+1$ and $j> \\frac{n}{i+1}$, then \n\\[\n[z^j] F_n(z) = [z^{j-1}] F_{n-i}(z) .\n\\]\n\\end{theorem}\n\\begin{proof}\nBy (\\ref{eq:coeff_F}), we need to show that\n\\begin{eqnarray*}\n&&\n \\sum_{s=ji }^{ n} p(n-s, < i) \\, \\# \\{ \\mu \\vdash s : \\textrm{all parts of } \\mu \\textrm{ are } \\geq i, \\sum_{a,d} a e_{\\mu(a,d)}=j \\}\n \\\\\n &&\n =\n \\sum_{s'= (j-1) i }^{n- i } p(n-s'-i, < i) \\,\n \\\\\n &&\n \\quad\n \\times \\quad\n \\# \\{ \\nu \\vdash s' : \\textrm{all parts of } \\nu \\textrm{ are } \\geq i, \\sum_{a,d} a e_{\\nu(a,d)}=j -1 \\}\n\\end{eqnarray*}\nThese two coefficients are equal provided we construct a bijection $T$ between\n\\begin{eqnarray*}\nU_s &=& \\{ \\mu \\vdash s : \\textrm{all parts of } \\mu \\textrm{ are } \\geq i, \\sum_{a,d} a e_{\\mu(a,d)}=j \\}\n \\\\\n V_{s'}\n &=&\n \\{ \\nu \\vdash s' : \\textrm{all parts of } \\nu \\textrm{ are } \\geq i, \\sum_{a,d} a e_{\\nu(a,d)}=j -1 \\}.\n\\end{eqnarray*}\nwhen $s' = s-i$. Let $\\mu \\in U_s$. By Lemma 7, the partition $\\mu$ must have a part equal to $i$. \nLet $\\nu=T(\\mu)$ be the partition of $s-i$ obtained by deleting one part from $\\mu$ of size $i$. It is easy\nto verify that $\\nu \\in V_{s'}$. The inverse of $T$ is simply adding a part of size $i$ to $\\nu \\in V_{s'}$.\n \\end{proof}\n \nIn other words, Theorem \\ref{thm:subsums} shows that if $m>i+1$ and $j> n\/(i+1)$, then the subsums satisfy\n$\\Lambda_{m,i}(n,j)= \\Lambda_{m,i}(n-i,j-1)$.\n\n\n\n\\section{Laurent Type Polynomials }\\label{section:laurent}\n\nA more general case consists of generating functions that involve $z$ raised to different powers; \nultimately, we might treat the case of the generating function \n\\[\nG(z,q) = \\prod_{(i,j) \\in \\mathbb{Z}^2 \\setminus (0,0)} {(1 - z^i q^j)}^{a_{ij}} \\, .\n\\]\n\n \nAn important example comes from the generating function for the crank statistic for partitions. \nLet\n\\[\nC(z,q) = \\prod_{k\\geq 1} \\frac{1-q^k}{(1-z q^k)(1-z^{-1} q^k)} = \\sum_{n=0}^\\infty \\, M_n(z) q^n \n\\]\nwhere $M_n(z)$ is a symmetric Laurent polynomial.\nFrom the definition of the crank, the coefficient of $z^{n-k}$ in $M_n(z)$,\nfor $ k \\leq \\frac{n}{2}$, equals the number of partitions of $k$ that include no 1s.\nIn particular, the coefficients of $M_n(z)$ stabilize in the ranges for powers $n-k$ and $-n+k$\nfor $0\\leq k \\leq \\lfloor n\/2 \\rfloor$. It is suggested in \\cite{BG2} that the zeros for the crank polynomial converge to\nthe unit circle. \nAnother example is the generating\nfunction\n\\[\n\\prod_{k=1}^\\infty \\frac{ 1}{ (1-z q^k)^{k}} \\, \\frac{1}{ (1-z^{-1} q^k)^k } \\, \\frac{1}{ (1-q^k)^{2k}}.\n\\]\nwhich comes from the Donaldson-Thomas Theory in algebraic geometry and whose asymptotics were studied in \\cite{siam10}.\n\n\n\n \n\\begin{lemma}\\label{thm:freezing1}\n Let $\\{A_n(z)\\}_{n=0}^\\infty$ be a sequence of polynomials, such that the degree of $A_n(z)$ is $n$, whose coefficients satisfy\n \\[\n [z^{n-k}] A_n(z) = [ z^{n+1-k}]A_{n+1}(z), \\quad 0 \\leq k \\leq n\/m\n \\]\n for some integer $m \\geq 2$.\n Let $\\{B_n(z)\\}_{n=0}^\\infty$ be another sequence of polynomials. Then the coefficients of the \n polynomial sequence $\\{F_n(z)\\}_{n=0}^\\infty$ \n where\n \\[\n F_n(z) = \\sum_{\\ell=0}^n \\, A_{\\ell}(z) B_{n-\\ell}(z^{-1})\n \\]\n also satisfy\n \\[\n [z^{n-k}] F_n(z) = [ z^{n+1-k}]F_{n+1}(z), \\quad 0 \\leq k \\leq n\/m.\n \\]\n \\end{lemma}\n \\begin{proof}\n Let $0 \\leq k \\leq \\lfloor n\/m \\rfloor$. Consider $[z^{n-k}] F_n(z)$. We have\n \\begin{eqnarray}\n [z^{n-k}] F_n(z)\n &=&\n [z^{n-k}] \\sum_{\\ell=0}^n A_{\\ell}(z) \\, B_{n-\\ell}(z^{-1})\n \\nonumber\n \\\\\n &=&\n\\sum_{\\ell=0}^n \\, \\sum_{a=0}^k \\, [ z^{ n-k+a}] A_\\ell(z) \\, [ z^{-a}] B_{n-\\ell}(z^{-1})\n \\nonumber\n \\\\\n &=&\n\\sum_{\\ell=n-k}^n \\, \\sum_{a=0}^{k} \\, [ z^{ n-k+a}] A_\\ell(z) \\, [ z^{-a}] B_{n-\\ell}(z^{-1})\n \\label{eq:first_sum}\n \\end{eqnarray}\nwhere in the last step we note that $[z^{ n-k+a}]A_\\ell(z)=0$ if $\\ell < n-k$.\n\nFor $[z^{n+1-k}]F_{n+1}(z)$, we reindex $\\ell$ by 1 and note that $[z^{n+1-k}]A_0(z)=0$ for all $n$, giving us the expression for $[z^{n+1-k}]F_{n+1}(z)$:\n \\begin{equation}\n \\label{eq:second_sum}\n\\sum_{\\ell=n-k}^{n} \\, \\sum_{a=0}^{k} \\, [ z^{ n+1-k+a}] A_{\\ell+1}(z) \\, [ z^{-a}] B_{n-\\ell}(z^{-1})\n\\end{equation}\nBy assumption, we know that\n\\[\n [ z^{ n-k+a}] A_\\ell(z) =[ z^{ n+1-k+a}] A_{\\ell+1}(z), \\quad n-k \\leq \\ell \\leq n\n\\]\nsince $0 \\leq k \\leq \\lfloor n\/m \\rfloor$. The coefficients of $[z^{-a}] B_{n-\\ell}(z^{-1})$ are the same, and consequently the two sums (\\ref{eq:first_sum}) and (\\ref{eq:second_sum}) are\nequal.\n \\end{proof}\n \n \n\\begin{lemma}\\label{lemma:convolution}\nGiven the generating function\n$G(z,q) = \\prod_{i \\geq 1} (1-zq^i)^{-a_i}$ where $a_1=1$, let $A_n(z) = [q^n] G(z,q)$. Let $Q(q) = \\sum_{j=0}^\\infty c_j q^j$. \nLet $F_n(z) = [q^n] G(z,q)Q(q)$. Then the tail coefficients of $F_n(z)$ stabilize; that is,\n\\[\n[z^{k}] F_n(z) = [z^{k+1}] F_{n+1}(z), \\quad k \\geq n\/2.\n\\]\n\\end{lemma}\n\\begin{proof}\nBy construction, the polynomials $F_n(z)$ have the form\n\\[\nF_n(z) = \\sum_{\\ell=0}^n c_{n-\\ell} A_\\ell(z).\n\\]\nLet $k \\geq n\/2$. Then the coefficients for $[z^{k}] F_n(z)$ and $ [z^{k+1}] F_{n+1}(z)$ are given by\n\\begin{eqnarray*}\n[z^{k}] F_n(z) = \\sum_{\\ell=k}^n c_{n-\\ell} [z^{k} ] A_\\ell(z),\n\\quad\n\\,[z^{k+1}]F_{n+1}(z) = \\sum_{j = k+1}^{n+1} c_{n+1-j} [z^{k+1}] A_{j}(z).\n\\end{eqnarray*}\nAs a consequence of Corollary 1, $c_{n-\\ell} [z^{k} ] A_\\ell(z) = c_{n+1-j} [z^{k+1}] A_{j}(z)$ for $j= \\ell+1$ and $k \\leq \\ell \\leq n$.\n\\end{proof}\n \n \\begin{theorem}\n If $a_1 = b_1 = 1$, and \n\\[\nG(z,q) = \\sum_{n=0}^\\infty F_n (z) q^n = \\prod_{i \\geq 1} {(1 - z q^i)}^{-a_i} {(1 - z^{-1} q^i)}^{-b_i} {(1 \\pm q^i)}^{c_i},\n\\]\nthen the coefficients of the Laurent polynomials $F_n(z)$ satisfy\n\\[\n[z^{n-k}] F_n(z) = [z^{n+1-k}] F_{n+1}(z), \\quad [z^{-(n-k)}] F_n(z) = [z^{-(n+1-k)}] F_{n+1}(z),\n\\]\nfor $0 \\leq k \\leq n\/2$.\n \\end{theorem}\n \n \n\n\n\\section{ Plane Overpartition Stabilization }\n\nA plane partition is an array of positive integers, conventionally justified to the upper left corner of the fourth quadrant, which are weakly descending left in rows and down in columns. A plane overpartition is a plane partition whose entries may be overlined or not according to\ncertain rules \\cite{CSV}: in each row, the last occurrence of an integer may be overlined\n(or not) and in every column, all but the first occurrence of an integer are overlined, while the first occurrence may or may not be overlined.\nIn \\cite[Proposition 4]{CSV}, the generating function for the weighted plane overpartitions is found to be \n\\[\n\\sum_{ \\Pi \\,\\,\\textrm{is a plane overpartition}}\nz^{ o(\\Pi)} q^{ | \\Pi |} =\n\\prod_{n=1}^\\infty \\frac{ (1+ zq^n)^n}{ (1-q^n)^{ \\lceil n\/2 \\rceil} (1-z^2q^n)^{ \\lfloor n\/2 \\rfloor } }\n\\]\nwhere $o( \\Pi)$ is the number of overlined parts of the plane\noverpartition $\\Pi$.\n\n\n\n\\begin{theorem}\nLet $G(z,q)$ be the generating function for the polynomials $F_n(z)$:\n\\[\nG(z,q)=\n\\prod_{n=1}^\\infty \\frac{ (1+ zq^n)^n}{ (1-q^n)^{ \\lceil n\/2 \\rceil} (1-z^2q^n)^{ \\lfloor n\/2 \\rfloor } }\n=\n\\sum_{n=0}^\\infty F_n(z) q^n .\n\\]\nThen the coefficients of the polynomials $F_n(z)$ satisfy the stabilization condition\n\\[\n[z^{k+1}] F_{n+1}(z) = [ z^k] F_n(z),\n\\]\nfor $k\\geq 2n\/3$.\n\\end{theorem}\n\\begin{proof}\nLet $\\{A_n(z)\\}$ be the polynomial sequence with generating function $G_A(z,q)$ where\n\\[\nG_A(z,q) = \\prod_{n=2}^\\infty \\frac{ 1}{ (1-z^2q^n)^{ \\lfloor n\/2 \\rfloor } } = \\sum_{n=0}^\\infty A_n(z) q^n\n\\]\nand $\\{ B_n(z)\\}$ with generating function $G_B(z,q)$:\n\\[\nG_B(z,q) = \\prod_{n=2}^\\infty (1+ zq^n)^n = \\sum_{n=0}^\\infty B_n(z) q^n .\n\\]\nEasily, we have that $\\deg(A_n) = 2 \\lfloor n\/2 \\rfloor$.\nBy Theorem \\ref{thm:stabilization2}, replacing $z$ by $z^2$, we also find\n\\[\n[z^k] A_n = [ z^{k+2}] A_{n+2}, \\quad k \\geq 2n\/3 .\n\\]\n\nThe degree of the polynomial $B_N(z)$ is the largest number of parts in a possible partition of $N$ drawn from a multiset of two 2s, three 3s, etc. For $N \\geq 21 = 2+2+3+3+3+4+4$, the average size of part is at least 3 and so \n\\[\n\\deg(B_N) \\leq \\frac{1}{3}N \\, \\text{.}\n\\]\n For smaller $N$, direct calculation shows that $\\deg(B_N) \\leq \\frac{2}{3}N$.\n\n\nIt is more convenient to work with the intermediate polynomials\n\\[\nQ_n(z) = [q^n] G_A(z,q) G_B(z,q) =\\sum_{\\ell=0}^n A_{n-\\ell}(z) \\, B_\\ell(z).\n\\]\n\nWe have the summation formula for $[z^k] Q_n(z)$:\n\\begin{eqnarray*}\n[z^k] Q_n(z)\n&=&\n\\sum_{\\ell=0}^n \\, \\sum_{a=0}^k [ z^{k-a}] A_{n - \\ell}(z) \\cdot [z^a] B_\\ell(z)\n\\\\\n\\\\\n&=&\n\\sum_{\\ell \\in I_{k,n} } \\, \\sum_{ a =0 }^{ \\deg(B_\\ell)} \\, [ z^{k-a}] A_{n - \\ell}(z) \\cdot [z^a] B_\\ell(z) \n\\end{eqnarray*}\nwhere $I_{k,n}$ is the set of indices \n\\[\nI_{k,n} =\n\\left\\{\n\\ell :\n\\deg(B_\\ell) + \\deg(A_{n - \\ell }) \\geq k, \\, 0 \\leq \\ell \\leq n\n\\right\\} .\n\\]\nWe have $I_{k+2, n+2} = I_{k,n} $ since for any $B_\\ell$ the matching $A_{n-\\ell}$ for sums of $k$ map to the matching $A_{n+2-\\ell}$ for sums of $k+2$.\n\n\nWe next show that\n\\[\n\\sum_{\\ell \\in I_{k,n} } \\, \\sum_{ a = 0}^{ \\deg(B_\\ell)} \\, [ z^{k-a}] A_{n - \\ell}(z) \\cdot [z^a] B_\\ell(z) \n=\n\\sum_{\\ell \\in I_{k+2,n+2} } \\, \\sum_{ a = 0}^{ \\deg(B_\\ell)} \\, [ z^{k+2-a}] A_{n+2 - \\ell}(z) \\cdot [z^a] B_\\ell(z) ;\n\\]\nTo do this, we need that all the terms $[ z^{k-a}] A_{n- \\ell}$ fall into the stable range of indices. \nNow\nto get into the stable range, we need\n\\[\nk-a \\geq \\frac{2}{3} ( n - \\ell) \n\\]\nwhere $0\\leq a \\leq \\deg B_\\ell$ and $\\ell \\in I_{k,n}$.\nWe consider the stronger condition\n\\begin{eqnarray*}\n\\frac{2}{3} n - \\deg B_\\ell &\\geq& \\frac{2}{3} ( n - \\ell) ,\n\\end{eqnarray*}\nwhich reduces to \n\\[\n\\deg(B_\\ell) \\leq \\frac{2}{3} \\ell, \\quad \\ell \\in I_{k,n} \n\\]\nwhich does indeed hold.\n\n\nOur next step is to define\n\\[\nP_n(z) = (1+zq) Q_n(z).\n\\]\n(Leaving out this factor earlier simplified the degree analysis, since without it we had the single exceptional case of $\\deg (B_1) = 1$.) We observe that\n\\begin{eqnarray*}\n\\,\n[z^k] P_n (z) &=& [ z^k] Q_n(z) + [ z^{k-1} ] Q_{n-1}(z),\n\\\\\n\\,\n[z^{k+1}] P_{n+1} (z) &=& [ z^{k+1}] Q_{n+1}(z) + [ z^{k}] Q_{n}(z) .\n\\end{eqnarray*}\nHence, we have\n\\[\n[z^k] P_n(z) = [z^{k+1}] P_{n+1}(z), \\quad k \\geq 2n\/3 .\n\\]\n\nTo finish the proof, we define a polynomial family $C_n(q)$ by\n\\[\n\\prod_{n=1}^\\infty \\frac{ 1 }{ (1-q^n)^{ \\lceil n\/2 \\rceil} } = \\sum_{n=0}^\\infty C_n(q) .\n\\]\nThen the polynomials $F_n(z)$ are given by\n\\[\nF_n(z) = [q^n] C(q) \\, \\sum_{\\ell =0}^\\infty P_\\ell (z) q^\\ell .\n\\]\nBy Lemma \\ref{lemma:convolution}, \nwe see that this construction maintains the stability of the coefficients of the polynomial family $\\{ P_\\ell(z)\\}$.\n\\end{proof}\n\n\\begin{corollary}\nLet $\\overline{pp}_k(n)$ be the number of plane overpartitions of $n$ with $k$ overlined parts. If $k \\geq 2n\/3$, then\n\\[\n\\overline{pp}_k(n) =\\overline{pp}_{k+1}(n+1).\n\\]\n\\end{corollary}\n \n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\\label{sec:form}\n\nIn multicellular organisms ranging from cellular amoebae to the human\nbody, it is essential for cells to communicate with each other. One\ncommon mechanism to initiate communication between cells that are not\nin close contact is for cells to secrete diffusible signaling\nmolecules into the extracellular space between the spatially\nsegregated units. Examples of this kind of signaling range from\ncolonies of the amoebae Dictyostelium discoideum, which release cAMP\ninto the medium where it diffuses and acts on each separate colony\n(cf.~\\cite{Goldbeter1990}), to some endocrine neurons that secrete a\nhormone to the extracellular medium where it influences the secretion\nof this hormone from a pool of such neurons\n(cf.~\\cite{Krasmanovic_etal2003}, \\cite{Li2008}), to the effect of\ncatalysts in surface science \\cite{riecke}, and to quorum sensing\nbehavior for various applications (cf.~\\cite{CLL}, \\cite{Muller1},\n\\cite{Muller2}). In many of these systems, the individual cells or\nlocalized units can exhibit sustained temporal oscillations. In this\nway, signaling through a diffusive chemical can often trigger\nsynchronous oscillations among all the units.\n\nIn this paper we provide a theoretical investigation of the mechanism\nthrough which this kind of synchronization occurs for a class of\ncoupled cell-bulk ODE-PDE models in bounded two-dimensional\ndomains. Our class of models consists of $m$ small cells with\nmulti-component intracellular dynamics that are coupled together by a\ndiffusion field that undergoes constant bulk decay. We assume that the\ncells can release a specific signaling molecule into the bulk region\nexterior to the cells, and that this secretion is regulated by both\nthe extracellular concentration of the molecule together with its\nnumber density inside the cells. Our aim is to characterize conditions\nfor which the release of the signaling molecule leads to the\ntriggering of some collective synchronous oscillatory behavior among\nthe localized cells. Our modeling framework is closely related to the\nstudy of quorum sensing behavior in bacteria done in \\cite{Muller1}\nand \\cite{Muller2} through the formulation and analysis of similar\ncoupled cell-bulk models in ${\\mathbb{R}}^3$. For this 3-D case, in\n\\cite{Muller1} and \\cite{Muller2} steady-state solutions were\nconstructed and large-scale dynamics studied in the case where the\nsignaling compartments have small radius of order ${\\mathcal\n O}(\\epsilon)$. However, due to the rapid ${1\/r}$ decay of the\nfree-space Green's function for the Laplacian in 3-D, it was shown in\n\\cite{Muller1} and \\cite{Muller2} that the release of the signaling\nmolecule leads to only a rather weak communication between the cells\nof the same ${\\mathcal O}(\\epsilon)$ order of the cell radius. As a\nresult, small cells in 3-D are primarily influenced by their own\nsignal, and hence no robust mechanism to trigger collective\nsynchronous oscillations in the cells due to Hopf bifurcations was\nobserved in \\cite{Muller1} and \\cite{Muller2}. We emphasize that the\nmodels of \\cite{Muller1} and \\cite{Muller2} are based on postulating a\ndiffusive coupling mechanism between distinct, spatially segregated,\ndynamically active sites. Other approaches for studying quorum sensing\nbehavior, such as in \\cite{NSSGM2014}, are based on reaction-diffusion\n(RD) systems, which adopt a homogenization theory approach to treat\nlarge populations or colonies of individual cells as a continuum\ndensity, rather than as discrete units as in \\cite{Muller1} and\n\\cite{Muller2}.\n\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.65\\textwidth,height=5.0cm]{twoDdiag.eps}\n\\caption{Schematic diagram showing the intracellular reactions and\n external bulk diffusion of the signal. The small shaded regions are\n the signaling compartments or cells.}\\label{fig:schem}\n\\end{center}\n\\end{figure}\n\nBefore discussing our main results, we first formulate and\nnon-dimensionalize our coupled cell-bulk model assuming that there is\nonly one signaling compartment $\\Omega_0$ inside the two-dimensional\ndomain $\\Omega$. The corresponding dimensionless model for $m>1$ small\nsignaling cells is given in (\\ref{mainbd}). Fig.~\\ref{fig:schem}\nshows a schematic plot of the geometry for $m$ small cells. We assume\nthat the cell can release a specific signaling molecule into the bulk\nregion exterior to the cell, and that this secretion is regulated by\nboth the extracellular concentration of the molecule together with its\nnumber density inside the cell. If ${\\mathcal U}({\\vecb X},T)$\nrepresents the concentration of the signaling molecule in the bulk\nregion $\\Omega\\backslash\\Omega_0$, then its spatial-temporal evolution\nin this region is assumed to be governed by the PDE model\n\\begin{subequations}\\label{f:mainbd}\n\\begin{equation}\\label{f:mainU}\n\\begin{aligned}\n {\\mathcal U}_T &= D_B\\Delta_{\\vecb X} {\\mathcal U} - k_B {\\mathcal U}\\,,\n \\qquad \\vecb X \\in \\Omega\\backslash\\Omega_{0}\\,; \\qquad\n\\partial_{n_{\\vecb X}} {\\mathcal U}= 0\\,,\\qquad \\vecb X \\in\\partial \\Omega\\,,\\\\\n D_B \\partial_{n_{\\vecb X}} {\\mathcal U}&= \\beta_1 {\\mathcal U} -\\beta_2 {\\mu}^1\n \\,, \\;\\;\\quad \\vecb X\\in\\partial \\Omega_{0}\\,, \\\\\n\\end{aligned}\n\\end{equation}\nwhere, for simplicity, we assume that the signaling compartment\n$\\Omega_0\\in\\Omega$ is a disk of radius $\\sigma$ centered at some\n$\\vecb X_0\\in \\Omega$. Inside the cell we assume that there are $n$\ninteracting species $\\vecb \\mu \\equiv (\\mu^1,\\ldots,\\mu^n)^T$ whose\ndynamics are governed by $n$-ODEs, with a source term representing the\nexchange of material across the cell membrane $\\partial\\Omega_0$, of\nthe form\n\\begin{equation}\\label{f:mainuj}\n\\frac{d \\vecb \\mu}{dT} = k_R \\mu_c \\vecb F\\left({{\\vecb \\mu}\/\\mu_c}\\right)\n+ \\vecb e_1 \\int_{\\partial \\Omega_{0}}\\left( \\beta_1 {\\mathcal U} -\n\\beta_2 \\mu^1\\right) \\, dS_{\\vecb X} \\,,\n\\end{equation}\n\\end{subequations}$\\!$ where $\\vecb e_1\\equiv (1,0,\\dots,0)^T$. Here $\\vecb \\mu$ is the\ntotal amount of the $n$ species inside the cell, while $k_R>0$ is the\nreaction rate for the dimensionless intracellular dynamics ${\\vecb\n F}(\\vecb u)$. The scalar $\\mu_c>0$ is a typical value for $\\vecb\n\\mu$.\n\nIn this coupled cell-bulk model, $D_B>0$ is the diffusion coefficient\nfor the bulk process, $k_B$ is the rate at which the signaling\nmolecule is degraded in the bulk, while $\\beta_1>0$ and $\\beta_2>0$\nare the dimensional influx (eflux) constants modeling the permeability\nof the cell wall. In addition, $\\partial_{n_{\\vecb X}}$ denotes either\nthe outer normal derivative of $\\Omega$, or the outer normal to\n$\\Omega_{0}$ (which points inside the bulk region). The flux\n$\\beta_1{\\mathcal U}-\\beta_2\\mu^1$ on the cell membrane models the\ninflux of the signaling molecule into the extracellular bulk region,\nwhich depends on both the external bulk concentration ${\\mathcal\n U}({\\vecb X},T)$ at the cell membrane $\\partial\\Omega_0$ as well as\non the intracellular concentration $\\mu^1$ within the cell. We assume\nthat only one of the intracellular species, $\\mu^1$, is capable of\nbeing transported across the cell membrane $\\partial\\Omega_0$ into the\nbulk. We remark that a related class of models was formulated and\nanalyzed in \\cite{bm1} and \\cite{bm2} in their study of the initiation\nof the biological cell cycle, where the dynamically active compartment\nis the nucleus of the biological cell.\n\nNext, we introduce our scaling assumption that the radius $\\sigma$ of\nthe cell is small compared to the radius of the domain, so that\n$\\epsilon\\equiv {\\sigma\/L}\\ll 1$, where $L$ is the length-scale of\n$\\Omega$. However, in order that the signaling compartment has a\nnon-negligible effect on the bulk process, we need to assume that\n$\\beta_1$ and $\\beta_2$ are both ${\\mathcal O}(\\epsilon^{-1})\\gg 1$ as\n$\\epsilon\\to 0$. In this way, in Appendix \\ref{app:A} we show that\n(\\ref{f:mainbd}) reduces to the dimensionless coupled system\n\\begin{subequations}\\label{f4:mainbd}\n\\begin{equation}\\label{f4:mainU}\n\\begin{aligned}\n \\tau { U}_t &= D \\Delta_{\\vecb x} {U} - {U}\\,, \\qquad \\vecb x \\in\n {\\Omega}\\backslash{\\Omega}_{0}\\,; \\qquad\n \\partial_{n_{\\vecb x}} {U} = 0\\,,\\qquad \\vecb x \\in\\partial\n {\\Omega}\\,,\\\\ \\epsilon D \\partial_{n_{\\vecb x}} {U}&= d_1 { U}\n - d_2 u^1 \\,, \\;\\;\\quad \\vecb x\\in\\partial {\\Omega}_{0}\\,, \\\\\n\\end{aligned}\n\\end{equation}\nwhere ${\\Omega}_0$ is a disk of radius $\\epsilon\\ll 1$ centered\nat some $x_0\\in {\\Omega}$. The bulk process is coupled to the\nintracellular dynamics by\n\\begin{equation}\\label{f4:mainuj}\n \\frac{d \\vecb u}{dt} = \\vecb F\\left(\\vecb u\\right)\n+ \\frac{\\vecb e_1}{\\tau \\epsilon} \\int_{\\partial {\\Omega}_{0}}\n \\left( d_1 {U} - d_2 u^1\\right) \\, dS_{\\vecb x} \\,.\n\\end{equation}\n\\end{subequations}$\\!$\nThe four ${\\mathcal O}(1)$ dimensionless parameters in\n(\\ref{f4:mainbd}) are $\\tau$, $D$, $d_1$, and $d_2$, defined by\n\\begin{equation}\\label{form:par}\n \\tau \\equiv \\frac{k_R}{k_B} \\,, \\qquad D \\equiv \\frac{D_B}{k_B L^2} \\,,\n\\qquad \\beta_1 \\equiv \\left(k_B L\\right) \\frac{d_1}{\\epsilon} \\,, \\qquad\n \\beta_2 \\equiv \\left( \\frac{k_B}{L} \\right) \\frac{d_2}{\\epsilon} \\,.\n\\end{equation}\nWe remark that the limit $\\tau\\ll 1$ ($\\tau\\gg 1$) corresponds to when\nthe intracellular dynamics is very slow (fast) with respect to the\ntime-scale of degradation of the signaling molecule in the bulk. The\nlimit $D\\gg 1$ corresponds to when the bulk diffusion length\n$\\sqrt{D_B\/k_B}$ is large compared to the length-scale $L$ of the confining\ndomain. \n\nA related class of coupled cell-bulk models involving two bulk\ndiffusing species, and where the cells, centered at specific spatial\nsites, are modeled solely by nonlinear flux boundary conditions, have\nbeen used to model cellular signal cascades (cf.~\\cite{levy1},\n\\cite{levy2}), and the effect of catalyst particles on chemically\nactive substrates (cf.~\\cite{Peirce}, \\cite{Glass}). In contrast to\nthese nonlinear flux-based models, in the coupled cell-bulk models of\n\\cite{Muller1} and \\cite{Muller2}, and the one considered herein, the\ncells are not quasi-static but are, instead, dynamically active\nunits. Our main goal for (\\ref{f4:mainbd}) and (\\ref{mainbd}) is to\ndetermine conditions that lead to the triggering of synchronized\noscillations between the dynamically active cells. Related 1-D\ncell-bulk models, where the cells are dynamically active units at the\nends of a 1-D spatial domain, have been analyzed in\n\\cite{Gomez-Marin2007}--\\cite{glnw}.\n\nOur analysis of the 2-D coupled cell-bulk model (\\ref{f4:mainbd}), and\nits multi-cell counterpart (\\ref{mainbd}), which extends the 3-D\nmodeling paradigm of \\cite{Muller1} and \\cite{Muller2}, has the\npotential of providing a theoretical framework to model quorum sensing\nbehavior in experiments performed in petri dishes, where cells live on\na thin substrate. In contrast to the assumption of only one active\nintracellular component used in \\cite{Muller1} and \\cite{Muller2}, in\nour study we will allow for $m$ small spatially segregated cells with\nmulti-component intracellular dynamics. We will show for our 2-D case\nthat the communication between small cells through the diffusive\nmedium is much stronger than in 3-D, and leads in certain parameter\nregimes to the triggering of synchronous oscillations, which otherwise\nwould not be present in the absence of any cell-bulk coupling. In\naddition, when $D={\\mathcal O}(1)$, we show that the spatial\nconfiguration of small cells in the domain is an important factor in\ntriggering collective synchronous temporal instabilities in the cells.\n\nThe outline of this paper is as follows. In \\S \\ref{sec:bd} we use the\nmethod of matched asymptotic expansions to construct steady-state\nsolutions to our 2-D multi-cell-bulk model (\\ref{mainbd}), and we\nderive a globally coupled eigenvalue problem whose spectrum\ncharacterizes the stability properties of the steady-state. In our 2-D\nanalysis, the interaction between the cells is of order\n$\\nu\\equiv-1\/\\log\\epsilon$, where $\\epsilon\\ll 1$ is the assumed\ncommon radius of the small circular cells. In the distinguished limit\nwhere the bulk diffusion coefficient $D$ is of the asymptotic order\n$D=\\mathcal{O}(\\nu^{-1})$, in \\S \\ref{sec:largeD} we show that the\nleading order approximate steady-state solution and the associated\nlinear stability problem are both independent of the spatial\nconfigurations of cells and the shape of the domain. In this regime,\nwe then show that the steady-state solution can be destabilized by\neither a synchronous perturbation in the cells or by $m-1$ possible\nasynchronous modes of instability. In \\S \\ref{sec:largeD}\nleading-order-in-$\\nu$ limiting spectral problems when $D={D_0\/\\nu}$,\nwith $\\nu\\ll 1$, for both these classes of instabilities are\nderived. In \\S \\ref{sec:examples}, we illustrate our theory for\nvarious intracellular dynamics. When there is only a single\ndynamically active intracellular component, we show that no triggered\noscillations can occur. For two specific intracellular reaction\nkinetics involving two local species, modeled either by Sel'kov or\nFitzhugh-Nagumo (FN) dynamics, in \\S \\ref{sec:examples} we perform\ndetailed analysis to obtain Hopf bifurcation boundaries, corresponding\nto the onset of either synchronous or asynchronous oscillations, in\nvarious parameter planes. In addition to this detailed stability\nanalysis for the $D=\\mathcal{O}(\\nu^{-1})$ regime, in \\S\n\\ref{sec:odes} we show for the case of one cell that when $D\\gg\n\\mathcal{O}(\\nu^{-1})$ the coupled cell-bulk model is effectively\nwell-mixed and its solutions can be well-approximated by a\nfinite-dimensional system of nonlinear ODEs. The analytical and\nnumerical study of these limiting ODEs in \\S \\ref{sec:odes} reveals\nthat their steady-states can be destabilized through a Hopf\nbifurcation. Numerical bifurcation software is then used to show the\nexistence of globally stable time-periodic solution branches that are\nintrinsically due to the cell-bulk coupling. For the $D={\\mathcal\n O}(1)$ regime, where the spatial configuration of the cells in the\ndomain is an important factor, in \\S \\ref{sec:finite_d} we perform a\ndetailed stability analysis for a ring-shaped pattern of cells that is\nconcentric within the unit disk. For this simple spatial configuration\nof cells, phase diagrams in the $\\tau$ versus $D$ parameter space, for\nvarious ring radii, characterizing the existence of either synchronous\nor asynchronous oscillatory instabilities, are obtained for the case\nof Sel'kov intracellular dynamics. These phase diagrams show that\ntriggered synchronous oscillations can occur when cells become more\nspatially clustered. In \\S \\ref{sec:finite_d} we also provide a clear\nexample of quorom sensing behavior, characterized by the triggering of\ncollective dynamics only when the number of cells exceeds a critical\nthreshold. Finally, in \\S \\ref{sec:disc} we briefly summarize our main\nresults and discuss some open directions.\n\nOur analysis of synchronous and asynchronous instabilities for\n(\\ref{mainbd}) in the $D={\\mathcal O}(\\nu^{-1})$ regime, where the\nstability thresholds are to, to leading-order, independent of the\nspatial configuration of cells, has some similarities with the\nstability analysis of \\cite{wei1}, \\cite{wei2}, \\cite{rozada}, and\n\\cite{chen} (see also the references therein) for localized spot\nsolutions to various activator-inhibitor RD systems with short range\nactivation and long-range inhibition. In this RD context, when the\ninhibitor diffusivity is of the order ${\\mathcal O}(\\nu^{-1})$,\nlocalized spot patterns can be destabilized by either synchronous or\nasynchronous perturbations, with the stability thresholds being, to\nleading-order in $\\nu$, independent of the spatial configuration of\nthe spots in the domain. The qualitative reason for this similarity\nbetween the coupled cell-bulk and localized spot problems is\nintuitively rather clear. In the RD context, the inhibitor diffusion\nfield is the long-range ``bulk'' diffusion field, which mediates the\ninteraction between the ``dynamically active units'', consisting of\n$m$ spatially segregated localized regions of high activator\nconcentration, each of which is is self-activating. In this RD\ncontext, asynchronous instabilities lead to asymmetric spot patterns,\nwhile synchronous oscillatory instabilities lead to collective\ntemporal oscillations in the amplitudes of the localized spots\n(cf.~\\cite{wei1}, \\cite{wei2}, \\cite{rozada}, and \\cite{chen}). A more\ndetailed discussion of this analogy is given in Remark \\ref{remark}.\n\nFinally, we remark that the asymptotic framework for the construction\nof steady-state solutions to the cell-bulk model (\\ref{mainbd}), and\nthe analysis of their linear stability properties, relies heavily on\nthe methodology of strong localized perturbation theory\n(cf.~\\cite{ward}). Related problems where such techniques are used\ninclude \\cite{ktcw}, \\cite{pill}, \\cite{levy1}, and \\cite{levy2}.\n\n\\setcounter{equation}{0}\n\\setcounter{section}{1}\n\\section{Analysis of the Dimensionless 2-D Cell-Bulk System}\\label{sec:bd}\n\nWe first generalize the one-cell model of \\S \\ref{sec:form} by\nformulating a class of dimensionless coupled cell-bulk dynamics that\ncontains $m$ small, disjoint, cells or compartments that are scattered\ninside the bounded two-dimensional domain $\\Omega$. We assume that\neach cell is a small disk of a common radius $\\epsilon\\ll 1$ that\nshrinks to a point $\\vecb x_j\\in \\Omega$ as $\\epsilon\\to 0$, and that\nare well-separated in the sense that dist$(\\vecb x_i, \\vecb x_j)=O(1)$ for\n$i\\neq j$ and dist$(\\vecb x_j, \\partial \\Omega)=O(1)$ for\n$j=1,\\ldots,m$, as $\\epsilon\\to 0$.\n\nAs motivated by the dimensional reasoning provided in \\S\n\\ref{sec:form}, if $U({\\vecb x},t)$ is the dimensionless\nconcentration of the signaling molecule in the bulk region between the\ncells, then in this region it satisfies the dimensionless PDE \n\\begin{subequations}\\label{mainbd}\n\\begin{equation}\\label{mainU}\n\\begin{aligned}\n\\tau U_t &= D\\Delta U -U\\,,\\qquad \\vecb x\\in \\Omega\\backslash\n\\cup_{j=1}^m\\Omega_{\\epsilon_j}\\,; \\qquad\n\\partial_nU = 0\\,, \\qquad \\vecb x\\in\\partial \\Omega\\,,\\\\\n\\epsilon D \\partial_{n_j} U&=d_1U-d_2u^1_j\\,, \\;\\;\\quad \\vecb x\\in\\partial \n\\Omega_{\\epsilon_j}\\,, \\quad j=1,\\ldots,m\\,.\n\\end{aligned}\n\\end{equation}\nHere $D>0$ is the effective diffusivity of the bulk, $d_1>0$ and\n$d_2>0$ are the dimensionless influx (eflux) constants modeling the\npermeability of the cell membrane, $\\partial_n$ denotes the outer normal\nderivative of $\\Omega$, and $\\partial_{n_j}$ denotes the outer normal\nto each $\\Omega_{\\epsilon_j}$, which points inside the bulk\nregion. The signaling cell, or compartment,\n$\\Omega_{\\epsilon_j}$ is assumed to lie entirely within $\\Omega$. The\nflux $d_1U-d_2u^1_j$ on each cell membrane models the influx of the\nsignaling molecule into the extracellular bulk region, which depends\non both the external bulk concentration $U({\\vecb x},t)$ at the cell\nmembrane $\\partial\\Omega_{\\epsilon_j}$ as well as on the amount\n$u_j^1$ of one of the intracellular species within the $j$-th cell. We\nsuppose that inside each of the $m$ cells there are $n$ interacting\nspecies, with intracellular dynamics\n\\begin{equation}\\label{mainuj}\n\\frac{d \\vecb u_j}{dt} = \\vecb F_j(\\vecb u_j)+ \\frac{\\vecb e_1}{\\epsilon\n\\tau} \\int_{\\partial \\Omega_{\\epsilon_j}}( d_1 U -d_2 u^1_j)\\, ds \\,,\n\\end{equation}\n\\end{subequations}$\\!$ where $\\vecb e_1\\equiv (1,0,\\dots,0)^T$. Here $\\vecb u_j =\n(u^1_j,\\dots,u^n_j)^T$ is the mass of the $n$ species inside the\n$j$-th cell and $\\vecb F_j(\\vecb u_j)$ is the vector nonlinearity\nmodeling the reaction dynamics within the $j$-th cell. The integration\nin (\\ref{mainuj}) is over the boundary $\\partial\\Omega_{\\epsilon_j}$\nof the compartment. Since its perimeter has length $|\\partial\n\\Omega_{\\epsilon_j}|={\\mathcal O}(\\epsilon)$, this source term for the\nODE in (\\ref{mainuj}) is ${\\mathcal O}(1)$ as $\\epsilon\\to 0$.\nThe dimensionless parameters $D$, $\\tau$, $d_1$, and $d_2$, are related\nto their dimensional counterparts by (\\ref{form:par}). The qualitative\ninterpretation of the limits $\\tau\\ll 1$, $\\tau\\gg 1$, and $D\\gg 1$, were\ndiscussed following (\\ref{form:par}).\n\n\\subsection{The Steady-State Solution for the $m$ Cells System}\nWe construct a steady-state solution to (\\ref{mainbd}) under the\nassumption that the cells are well-separated in the sense described\npreceding (\\ref{mainbd}). In \\S \\ref{sec:stability} we will then\nformulate the linear stability problem for this steady-state solution.\n\nSince in an ${\\mathcal O}(\\epsilon)$ neighborhood near each cell the\nsolution $U$ has a sharp spatial gradient, we use the method of\nmatched asymptotic expansions to construct the steady-state solution\nto (\\ref{mainbd}). In the inner region near the $j$-th cell, we\nintroduce the local variables $U_j$ and $\\vecb y$, defined by $\\vecb y\n= \\epsilon^{-1}(\\vecb x-\\vecb x_j)$ and $U_j(\\vecb y) = U({\\vecb x_j}\n+ \\epsilon {\\vecb y}) \\,,$ so that (\\ref{mainU}) transforms to\n\\begin{equation}\\label{ssy}\nD\\Delta_{\\vecb y} U_j -\\epsilon^2 U_j= 0\\,,\\qquad |\\vecb y|>1\\,; \\qquad\nD\\partial_{n_j} U_j = d_1 U_j-d_2 u^1_j\\,, \\qquad |\\vecb y|=1\\,.\n\\end{equation}\nWe look for a radially symmetric solution to (\\ref{ssy}) in the form\n$U_j=U_j(\\rho)$, where $\\rho\\equiv |{\\vecb y}|$ and $\\Delta_{\\vecb\n y}=\\partial_{\\rho\\rho}+\\rho^{-1}\\partial_\\rho$ denotes the radially\nsymmetric part of the Laplacian. Therefore, to leading order, we have that\n$U_j(\\rho)$ satisfies\n\\begin{equation}\\label{main:ssrho}\n\\partial_{\\rho\\rho}U_j+\\rho^{-1}\\partial_\\rho U_j = 0\\,,\\qquad 1<\\rho<\\infty\\,;\n\\qquad D\\frac{\\partial U_j}{\\partial \\rho} = d_1 U_j-d_2 u^1_j\\,, \\qquad \n\\rho =1\\,.\n\\end{equation}\nThe solution to (\\ref{main:ssrho}) in terms of a constant\n$S_j$, referred to as the {\\em source strength} of the $j$-th cell, is\n\\begin{equation}\\label{chijs}\nU_j = S_j \\log\\rho +\\chi_j\\,, \\qquad \\chi_j = \n\\frac{1}{d_1}(DS_j+d_2u^1_j)\\,,\\qquad j = 1, \\dots, m\\,.\n\\end{equation} \nThe constant $S_j$ will be determined below upon matching the inner\nsolutions to the outer solution.\n\nFrom the steady-state of the intracellular dynamics (\\ref{mainuj})\ninside each cell, we find that the source strength $S_j$ and the\nsteady-state solution $\\vecb u_j$ satisfy the nonlinear algebraic\nsystem \n\\begin{equation}\\label{bduj}\n\\vecb F_j(\\vecb u_j)+ \\frac{2\\pi D}{ \\tau} S_j \\vecb e_1=0\\,.\n\\end{equation}\nIn principle, we can determine $u^1_j$ in terms of the unknown $S_j$\nas $u^1_j = u^1_j(S_j)$. The other values $u_2^j,\\dots, u_n^j$ also\ndepend on $S_j$. Next, in terms of $u_j^1$, we will derive a system of\nalgebraic equations for $S_1,\\ldots,S_m$, which is coupled to\n(\\ref{bduj}).\n\nUpon matching the far-field behavior of the inner solution\n(\\ref{chijs}) to the outer solution, we obtain the outer problem\n\\begin{equation}\\label{bdouterU}\n\\begin{aligned}\n \\Delta U-\\varphi_0^2 U &=0\\,, \\qquad \\vecb x\\in\\Omega\\backslash\n\\lbrace{\\vecb x_1,\\ldots,\\vecb x_m\\rbrace}\\,; \\qquad \\partial_n U=0\\,, \n \\qquad \\vecb x\\in\\partial \\Omega\\,,\\\\\n U &\\sim S_j\\log|\\vecb x-\\vecb x_j|+\\frac{S_j}{\\nu}+\\chi_j\\,, \\quad \n \\text{as }\\;\\vecb x\\rightarrow \\vecb x_j\\,, \\quad j=1,\\ldots,m\\,,\n\\end{aligned}\n\\end{equation}\nwhere we have defined $\\varphi_0$ and $\\nu\\ll 1$ by\n\\begin{equation}\n \\varphi_0\\equiv1\/\\sqrt{D}\\,,\\qquad \\nu\\equiv {-1\/\\log\\epsilon}\\,.\n\\end{equation}\nWe remark that the singularity condition in (\\ref{bdouterU}) is\nderived by matching the outer solution for $U$ to the far-field\nbehavior of the inner solution (\\ref{chijs}). We then introduce the\nreduced-wave Green's function $G(\\vecb x;\\vecb x_j)$ satisfying\n\\begin{subequations} \\label{bdgreenss_all}\n\\begin{equation}\\label{bdgreenss}\n\\Delta G-\\varphi_0^2 G = -\\delta(\\vecb x-\\vecb x_j)\\,,\\qquad \\vecb x\\in \n\\Omega \\,; \\qquad \\partial_n G = 0\\,, \\qquad \\vecb x\\in \\partial \\Omega\\,.\n\\end{equation}\nAs $\\vecb x\\rightarrow \\vecb x_j$, this Green's function has the local behavior\n\\begin{equation}\\label{bdgreenss_reg}\nG(\\vecb x;\\vecb x_j)\\sim -\\frac{1}{2\\pi}\\log|\\vecb x - \\vecb x_j|+R_j+o(1)\\,, \n\\qquad \\mbox{as} \\quad \\vecb x\\rightarrow \\vecb x_j\\,,\n\\end{equation}\n\\end{subequations}$\\!$\nwhere $R_j=R_j(\\vecb x_j)$ is called the regular part of $G(\\vecb\nx;\\vecb x_j)$ at $\\vecb x=\\vecb x_j$. In terms of $G(\\vecb x;\\vecb\nx_j)$, the solution to (\\ref{bdouterU}) is\n\\begin{equation}\\label{bdUso}\nU(\\vecb x)=-2\\pi\\sum_{i=1}^{m}S_iG(\\vecb x, \\vecb x_i)\\,.\n\\end{equation}\nBy expanding $U$ as $\\vecb x\\rightarrow \\vecb x_j$, and equating the\nresulting expression with the required singularity behavior in\n(\\ref{bdouterU}), we obtain the following algebraic system\nfor $\\vecb S=(S_1,\\ldots,S_m)^T$, which we write in matrix form as\n\\begin{equation}\\label{bdsysS2m}\n\\left(1+\\frac{D\\nu}{d_1}\\right)\\vecb S+2\\pi\\nu {\\mathcal G}\\vecb S=-\n\\frac{d_2}{d_1}\\nu \\vecb u^1\\,.\n\\end{equation}\nHere the Green's matrix ${\\mathcal G}$, with matrix entries\n$({\\mathcal G})_{ij}$, and the vector $\\vecb u^1$, whose $j$-th element\nis the first local species in the $j$-th cell, are given by\n\\begin{equation*}\n({\\mathcal G})_{ii} = R_{i} \\,, \\qquad ({\\mathcal G})_{ij} = \nG(\\vecb x_i;\\vecb x_j) \\equiv G_{ij} \\,, \\quad i\\neq j\\,; \\qquad \\vecb u^1\\equiv\n\\left(u_1^1,\\ldots,u_m^1\\right)^T \\,.\n\\end{equation*}\nSince $G_{ji}=G_{ij}$, by the reciprocity of the Green's function,\n${\\mathcal G}$ is a symmetric matrix.\n\nTogether with (\\ref{bduj}), (\\ref{bdsysS2m}) provides an approximate\nsteady-state solution for $\\vecb u$, which is coupled to the source\nstrengths $\\vecb S$. It is rather intractable analytically to write\ngeneral conditions on the nonlinear kinetics to ensure the existence\nof a solution to the coupled algebraic system (\\ref{bduj}) and\n(\\ref{bdsysS2m}). As such, in \\S \\ref{sec:examples} below we will\nanalyze in detail some specific choices for the nonlinear kinetics. We\nremark that even if we make the assumption that the nonlinear kinetics\nin the cells are identical, so that $\\vecb F_j=\\vecb F$ for\n$j=1,\\ldots,m$, we still have that $S_j$ and $\\vecb u^1$ depend on $j$\nthrough the Green's interaction matrix ${\\mathcal G}$, which depends\non the spatial configuration $\\lbrace{\\vecb x_1,\\ldots,\\vecb\n x_m\\rbrace}$ of the cells within $\\Omega$.\n\nIn summary, after solving the nonlinear algebraic system (\\ref{bduj})\nand (\\ref{bdsysS2m}), the approximate steady-state solution for $U$ is\ngiven by (\\ref{bdUso}) in the outer region, defined at ${\\mathcal\n O}(1)$ distances from the cells, and (\\ref{chijs}) in the\nneighborhood of each cell. This approximate steady-state solution is\naccurate to all orders in $\\nu$, since our analysis has effectively\nsummed an infinite order logarithmic expansion in powers of $\\nu$ for\nthe steady-state solution. Related 2-D problems where infinite\nlogarithmic expansions occur for various specific applications were\nanalyzed in \\cite{ktcw} and \\cite{pill} (see also the references therein).\n\n\\subsection{Formulation of the Linear Stability Problem}\\label{sec:stability}\nNext, we consider the linear stability of the steady-state solution\nconstructed in the previous subsection. We perturb this steady-state\nsolution, denoted here by $U_e(\\vecb x)$ in the bulk region and\n$\\vecb u_{e,j}$ in the $j$-th cell as $U=U_e + e^{\\lambda\n t}\\eta({\\vecb x})$ and $\\vecb u_j = \\vecb u_{e,j}+ e^{\\lambda\n t}\\vecb \\phi_j$. Upon substituting this perturbation into\n(\\ref{mainbd}), we obtain in the bulk region that\n\\begin{subequations}\\label{bdeigall}\n\\begin{equation}\\label{bdeigeta}\n\\begin{aligned}\n\\tau \\lambda \\eta &=D\\Delta \\eta-\\eta\\,,\\qquad \n \\vecb x\\in\\Omega\\backslash\\cup_{j=1}^m\\Omega_{\\epsilon_j}\\,; \\qquad\n\\partial_n \\eta = 0\\,,\\qquad \\vecb x\\in\\partial \\Omega\\,,\\\\\n\\epsilon D\\partial_{n_j}\\eta &= d_1\\eta-d_2\\vecb \\phi_j^1\\,,\n\\qquad \\vecb x\\in\\partial \\Omega_{\\epsilon_j}\\,, \\qquad j=1,\\ldots,m\\,.\n\\end{aligned}\n\\end{equation}\nWithin the $j$-th cell the linearized problem is\n\\begin{equation}\\label{bdeigphij}\n\\lambda \\phi_j=J_j\\phi_j+ \\frac{\\vecb e_1}{\\epsilon \\tau} \n \\int_{\\partial\\Omega_{\\epsilon_j}}\\left(d_1\\eta-d_2\\vecb \\phi_j^1\\right) \\, ds\\,,\n\\end{equation}\n\\end{subequations}$\\!$ \nwhere $J_j$ denotes the Jacobian matrix of the nonlinear\nkinetics $\\vecb F_j$ evaluated at $\\vecb u_{e,j}$. We now study\n(\\ref{bdeigall}) in the limit $\\epsilon\\to 0$ using the method of\nmatched asymptotic expansions. The analysis will provide a limiting\nglobally coupled eigenvalue problem for $\\lambda$, from which we can\ninvestigate possible instabilities.\n\nIn the inner region near the $j$-th cell, we introduce the local\nvariables $\\vecb y=\\epsilon^{-1}(\\vecb x-\\vecb x_j)$, with\n$\\rho=|{\\vecb y}|$, and let $\\eta_j(\\vecb y)=\\eta(\\vecb x_j + \\epsilon\n\\vecb y)$. We will look for the radially symmetric\neigenfunction $\\eta_j$ in the inner variable $\\rho$. Then, from\n(\\ref{bdeigeta}), upon neglecting higher order algebraic terms in $\\epsilon$,\nthe leading order inner problem becomes\n\\begin{equation}\\label{bdeig}\n\\partial_{\\rho\\rho}\\eta_j+\\rho^{-1}\\partial_\\rho\\eta_j= 0\\,, \\qquad \n1<\\rho<\\infty\\,; \\qquad\nD\\frac{\\partial \\eta_j}{\\partial \\rho} = d_1\\eta_j-d_2\\phi_j^1\\,,\n\\qquad \\rho=1\\,,\n\\end{equation}\nwhich has the solution\n\\begin{equation}\\label{bdBj}\n\\eta_j=c_j\\log \\rho+B_j\\,, \\qquad B_j=\\frac{1}{d_1}(Dc_j+d_2\\phi_j^1)\\,,\n\\end{equation}\nwhere $c_j$ is an unknown constant to be determined. Then, upon\nsubstituting (\\ref{bdBj}) into (\\ref{bdeigphij}), we obtain that\n\\begin{equation}\\label{bdeig1}\n(J_j-\\lambda I)\\vecb \\phi_j+ \\frac{2\\pi D}{ \\tau} c_j\\vecb\n e_1=0\\,, \\qquad j=1,\\ldots, m\\,.\n\\end{equation}\n\nIn the outer region, defined at ${\\mathcal O}(1)$ distances from the cells,\nthe outer problem for the eigenfunction $\\eta(\\vecb x)$ is\n\\begin{equation}\\label{bdeigout}\n\\begin{aligned}\n &\\Delta \\eta-\\frac{(1+\\tau\\lambda)}{D}\\eta =0, \\qquad \\vecb x\\in \n\\Omega\\backslash\\lbrace{\\vecb x_1,\\ldots,\\vecb x_m\\rbrace} \\,; \\qquad\n \\partial_n \\eta =0,\\qquad \\vecb x\\in \\partial \\Omega\\,,\\\\\n &\\eta \\sim c_j\\log|\\vecb x-\\vecb x_j|+\\frac{c_j}{\\nu}+B_j, \\qquad \n \\mbox{as} \\quad \\vecb x\\rightarrow \\vecb x_j\\,, \\quad j=1,\\ldots, m\\,,\n\\end{aligned}\n\\end{equation}\nwhere $\\nu\\equiv -{1\/\\log\\epsilon}$. We remark that the singularity\ncondition in (\\ref{bdeigout}) as $\\vecb x\\rightarrow \\vecb x_j$ is\nderived by matching the outer solution for $\\eta$ to the far-field\nbehavior of the inner solution (\\ref{bdBj}). To solve\n(\\ref{bdeigout}), we introduce the eigenvalue-dependent Green's\nfunction $G_{\\lambda}(\\vecb x; \\vecb x_j)$, which satisfies\n\\begin{equation}\\label{bdgneig}\n\\begin{aligned}\n&\\Delta G_{\\lambda}-\\varphi_{\\lambda}^2G_{\\lambda}=-\\delta(\\vecb x-\\vecb x_j), \n\\qquad \\vecb x\\in\\Omega\\,; \\qquad \\partial_n G_\\lambda = 0\\,,\\qquad\n \\vecb x\\in\\partial \\Omega\\,,\\\\\n&G_\\lambda(\\vecb x;\\vecb x_j) \\sim -\n\\frac{1}{2\\pi}\\log|\\vecb x-\\vecb x_j|+R_{\\lambda,j}+ o(1)\\,, \\qquad \n\\mbox{as} \\quad \\vecb x\\rightarrow\\vecb x_j\\,,\n\\end{aligned}\n\\end{equation}\nwhere $R_{\\lambda, j}\\equiv R_{\\lambda}(\\vecb x_j)$ is the regular part of\n$G_\\lambda$ at $\\vecb=\\vecb x_j$. Here we have defined $\\varphi_\\lambda$ by\n\\begin{equation}\n\\varphi_\\lambda\\equiv\\sqrt{\\frac{1+\\tau\\lambda}{D}}\\,.\n\\end{equation}\nWe must choose the principal branch of $\\varphi_\\lambda$, which\nensures that $\\varphi_\\lambda$ is analytic in $\\mbox{Re}(\\lambda)>0$.\nFor the case of an asymptotically large domain $\\Omega$, this choice\nfor the branch cut, for which $\\mbox{Re}(\\theta_\\lambda)>0$, also\nensures that $G_\\lambda$ decays as $|\\vecb x-\\vecb x_j|\\to \\infty$.\n\nIn terms of $G_\\lambda(\\vecb x; \\vecb x_j)$, we can represent the\nouter solution $\\eta(\\vecb x)$ satisfying (\\ref{bdeigout}), as\n\\begin{equation}\\label{eig:out}\n\\eta(\\vecb x) = -2\\pi \\sum_{i=1}^mc_iG_\\lambda(\\vecb x,\\vecb x_i)\\,.\n\\end{equation}\nBy matching the singularity condition at $\\vecb x\\rightarrow\\vecb\nx_j$, we obtain a system of equations for $c_j$ as\n\\begin{equation}\\label{bdcj1}\n c_j + \\nu B_j=-2\\pi\\nu\\left(c_j R_{\\lambda_j}+\\sum_{i\\neq j}^{m} c_i \nG_{\\lambda, ij} \\right) \\,, \\qquad j=1,\\ldots,m \\,,\n\\end{equation}\nwhere $G_{\\lambda, ij }\\equiv G_\\lambda(\\vecb x_j;\\vecb x_i)$.\nUpon recalling that $B_j = \\frac{1}{d_1}(Dc_j+d_2\\phi_j^1)$ from (\\ref{bdBj}),\nwe can rewrite (\\ref{bdcj1}) in matrix form in\nterms of $\\vecb c =(c_1,\\ldots,c_m)^T$ as \n\\begin{equation}\\label{bdeig2m}\n\\left(1+\\frac{D\\nu}{d_1}\\right)\\vecb c+\\frac{d_2}{d_1}\\nu\\vecb\\phi^1+2\\pi \\nu\n {\\mathcal G}_\\lambda \\vecb c=0\\,.\n\\end{equation}\nHere we have defined the symmetric Green's matrix ${\\mathcal\n G}_\\lambda$, with matrix entries $({\\mathcal G})_{\\lambda,ij}$, and\nthe vector $\\vecb \\phi^1$ by\n\\begin{equation}\\label{gcep:green_red}\n({\\mathcal G})_{\\lambda,ii} = R_{\\lambda,i} \\,, \\qquad ({\\mathcal G})_{\\lambda, ij} = \nG_{\\lambda}(\\vecb x_i;\\vecb x_j) \\equiv G_{\\lambda, ij} \\,, \\quad i\\neq j\\,; \\qquad \n\\vecb \\phi^1\\equiv \\left(\\phi_1^1,\\ldots,\\phi_m^1\\right)^T \\,.\n\\end{equation}\nThe $j$-th entry of the vector $\\phi^1=(\\phi^1_1,\\cdots,\\phi^1_m)^T$\nis simply the first element in the eigenvector for the $j$-th cell.\nTogether with (\\ref{bdeig1}), the system (\\ref{bdeig2m}) will yield an\neigenvalue problem for $\\lambda$ with eigenvector $\\vecb c$.\n\nNext, we calculate $\\phi_1$ in terms of $\\vecb c$ from (\\ref{bdeig1})\nand then substitute the resulting expression into (\\ref{bdeig2m}). If\n$\\lambda$ is not an eigenvalue of $J_j$, (\\ref{bdeig1}) yields that\n$\\vecb \\phi_j = 2\\pi D\\tau^{-1} (\\lambda I-J_j)^{-1} c_j\\vecb\ne_1$. Upon taking the dot product with the $n$-vector $\\vecb\ne_1=(1,0,\\ldots,0)^T$, we get $\\phi_j^1 = 2\\pi D \\tau^{-1} c_j {\\vecb\n e_1}^T (\\lambda I-J_j)^{-1} {\\vecb e_1}$, which yields in vector form that\n\\begin{subequations} \\label{kdef}\n\\begin{equation}\\label{kdef:1}\n \\phi^1 = \\frac{2\\pi D}{ \\tau} {\\mathcal K} \\vecb c \\,,\n\\end{equation}\nwhere ${\\mathcal K}={\\mathcal K}(\\lambda)$ is the $m\\times m$ diagonal\nmatrix with diagonal entries\n\\begin{equation}\\label{kdef:2}\n {\\mathcal K}_j = {\\vecb e_1}^T (\\lambda I-J_j)^{-1} {\\vecb e_1} =\n \\frac{1}{\\det(\\lambda I - J_j)} \\vecb e_1 M_j^T \\vecb e_1 = \n \\frac{M_{j,11}}{\\det(\\lambda I - J_j)} \\,. \n\\end{equation}\n\\end{subequations}$\\!$\nHere $M_j$ is the $n\\times n$ matrix of cofactors of the matrix\n$\\lambda I-J_j$, with $M_{j,11}$ denoting the matrix entry in the first row\nand first column of $M_{j}$, given explicitly by\n\\begin{equation}\\label{kdef:m11}\nM_{j,11}=M_{j,11}(\\lambda) \n \\equiv \\det\\left(\n\\begin{array}{ccc}\n\\lambda - \n\\frac{\\partial {F}_j^2}{\\partial u_2}\\Big|_{\\vecb{u}=\\vecb{u}_{e,j}},\n&\\cdots, &-\\frac{\\partial {F}_j^2}{\\partial u_n}\n\\Big|_{\\vecb{u}=\\vecb{u}_{e,j}}\\\\\n\\cdots,&\\cdots,&\\cdots\\\\\n-\\frac{\\partial {F}_j^n}{\\partial u_2}\\Big|_{\\vecb{u}=\\vecb{u}_{e,j}}\n ,&\\cdots, &\\lambda \n-\\frac{\\partial {F}_j^n}{\\partial u_n}\\Big|_{\\vecb{u}=\\vecb{u}_{e,j}}\n\\end{array}\n\\right) \\,.\n\\end{equation}\nHere ${F}_j^2,\\ldots,{F}_j^n$ denote the components of\nthe vector $\\vecb F_j \\equiv (F_{j}^1,\\ldots,F_j^n)^T$, characterizing\nthe intracellular kinetics.\n\nNext, upon substituting (\\ref{kdef:1}) into (\\ref{bdeig2m}), we obtain the\nhomogeneous $m\\times m$ linear system\n\\begin{subequations} \\label{gcep:full}\n\\begin{equation}\n {\\mathcal M} \\vecb c = \\vecb 0 \\,, \n\\end{equation}\nwhere the $m\\times m$ matrix ${\\mathcal M}={\\mathcal M}(\\lambda)$ is\ndefined by\n\\begin{equation}\\label{gcep:mdef}\n\\mathcal{M}\\equiv \\left(1+\\frac{D\\nu}{d_1}\\right)I+2\\pi \\nu \\frac{d_2}\n {d_1 { \\tau}} D {\\mathcal K} + 2\\pi \\nu {\\mathcal\n G}_\\lambda\\,.\n\\end{equation}\n\\end{subequations}$\\!$ In (\\ref{gcep:mdef}), the diagonal matrix ${\\mathcal K}$ has\ndiagonal entries (\\ref{kdef:2}), and ${\\mathcal G}_\\lambda$ is the\nGreen's interaction matrix defined in (\\ref{gcep:green_red}), which\ndepends on $\\lambda$ as well as on the spatial configuration\n$\\lbrace{\\vecb x_1,\\ldots,\\vecb x_m \\rbrace}$ of the centers of the\nsmall cells within $\\Omega$.\n\nWe refer to (\\ref{gcep:full}) as the globally coupled eigenvalue\nproblem (GCEP). In the limit $\\epsilon\\to 0$, we conclude that $\\lambda$ is\na discrete eigenvalue of the linearized problem (\\ref{bdeigall}) if\nand only if $\\lambda$ is a root of the transcendental equation\n\\begin{equation}\\label{bdMdet}\n\\det\\mathcal{M}=0\\,.\n\\end{equation}\nTo determine the region of stability, we must seek conditions to\nensure that all such eigenvalues satisfy $\\mbox{Re}(\\lambda)<0$. The\ncorresponding eigenvector $\\vecb c$ of (\\ref{gcep:full}) gives the\nspatial information for the eigenfunction in the bulk via\n(\\ref{eig:out}).\n\nWe now make some remarks on the form of the GCEP. We first observe\nfrom (\\ref{gcep:mdef}) that when $D={\\mathcal O}(1)$, then to\nleading-order in $\\nu\\ll 1$, we have that ${\\mathcal M}\\sim I +\n{\\mathcal O}(\\nu)$. As such, when $D={\\mathcal O}(1)$, we conclude\nthat to leading order in $\\nu$ there are no discrete eigenvalues of\nthe linearized problem with $\\lambda={\\mathcal O}(1)$, and hence no\n${\\mathcal O}(1)$ time-scale instabilities. However, since\n$\\nu={-1\/\\log\\epsilon}$ is not very small unless $\\epsilon$ is\nextremely small, this prediction of no instability in the $D={\\mathcal\n O}(1)$ regime may be somewhat misleading at small finite\n$\\epsilon$. In \\S \\ref{sec:finite_d} we determine the roots of\n(\\ref{bdMdet}) numerically, without first assuming that $\\nu\\ll 1$,\nfor a ring-shaped pattern of cells within the unit disk $\\Omega$, for\nwhich the Green's matrix is cyclic. In the next section we will\nconsider the distinguished limit $D={\\mathcal O}(\\nu^{-1})\\gg 1$ for\n(\\ref{gcep:mdef}) where the linearized stability problem becomes\nhighly tractable analytically.\n\n\\setcounter{equation}{0}\n\\setcounter{section}{2}\n\\section{The Distinguished Limit of $D={\\mathcal O}(\\nu^{-1})\\gg 1$}\n\\label{sec:largeD}\n\nIn the previous section, we constructed the steady-state solution for\nthe coupled cell-bulk system (\\ref{mainbd}) in the limit $\\epsilon\\to\n0$ and we derived the spectral problem that characterizes the linear\nstability of this solution. In this section, we consider the\ndistinguished limit where the signaling molecule in the bulk diffuses\nrapidly, so that $D\\gg 1$. More specifically, we will consider the\ndistinguished limit where $D={\\mathcal O}(\\nu^{-1})$, and hence for\nsome $D_0={\\mathcal O}(1)$, we set\n\\begin{equation}\\label{bdD0}\nD={D_0\/\\nu} \\,.\n\\end{equation}\n\nFor $D={\\mathcal O}(\\nu^{-1})$, we determine a leading order\napproximation for the steady-state solution and the associated\nspectral problem. To do so, we first approximate the reduced-wave\nGreen's function for large $D$ by writing (\\ref{bdgreenss}) as\n\\begin{equation}\\label{bdgreenssD}\n\\Delta G-\\frac{\\nu}{D_0}G=-\\delta(\\vecb x-\\vecb x_j)\\,,\\qquad \\vecb \nx\\in\\Omega\\,; \\qquad \\partial_n G=0, \\qquad \\vecb x\\in\\partial\\Omega\\,.\n\\end{equation}\nThis problem has no solution when $\\nu=0$. Therefore, we expand \n$G=G(\\vecb x; \\vecb x_j)$ for $D={D_0\/\\nu} \\gg 1$ as\n\\begin{equation}\\label{bdgexp}\nG=\\frac{1}{\\nu}G_{-1}+G_0+\\nu G_1+\\ldots \\,.\n\\end{equation}\nUpon substituting (\\ref{bdgexp}) into (\\ref{bdgreenssD}), we equate\npowers of $\\nu$ to obtain a sequence of problems for $G_i$ for\n$i=-1,0,1$. This leads to the following two-term expansion for \n$G(\\vecb x;\\vecb x_j)$ and its regular part $R_j$ in the limit\n$D={D_0\/\\nu} \\gg 1$:\n\\begin{equation}\\label{bdgexp1}\nG(\\vecb x; \\vecb x_j)=\\frac{D_0}{\\nu |{\\Omega}|}+\nG_{0}(\\vecb x; \\vecb x_j)+\\cdots\\,, \\qquad R_j =\\frac{D_0}{\\nu |{\\Omega}|}+\nR_{0,j} + \\cdots \\,.\n\\end{equation}\nHere $G_{0}(\\vecb x; \\vecb x_j)$, with regular part $R_{0j}$, is the\nNeumann Green's function defined as the unique solution to\n\\begin{equation}\\label{g0:neum}\n\\begin{aligned}\n\\Delta G_{0}&=\\frac{1}{|{\\Omega}|}-\\delta(\\vecb x-\\vecb x_j)\\,, \\qquad \n\\vecb x\\in\\Omega\\,; \\qquad \\partial_n G_0=0\\,,\\quad\\vecb x\\in \n\\partial\\Omega\\,; \\qquad \\int_{\\Omega} G_0\\: d\\vecb x =0\\,,\\\\\nG_0(\\vecb x;\\vecb x_j) &\\sim -\\frac{1}{2\\pi}\\log|{\\vecb x-\\vecb x_j}|+\n R_{0,j}, \\qquad \\vecb x\\rightarrow\\vecb x_j\\,.\n\\end{aligned}\n\\end{equation}\n\nWe then substitute the expansion (\\ref{bdgexp1}) and $D={D_0\/\\nu}$ into\nthe nonlinear algebraic system (\\ref{bduj}) and (\\ref{bdsysS2m}), which\ncharacterizes the steady-state solution, to obtain that\n\\begin{equation}\\label{bdsysS2m1}\n\\left(1+\\frac{D_0}{d_1}\\right)\\vecb S+\\frac{2\\pi m D_0}{|\\Omega|} \n {\\mathcal E} \\vecb S + 2\\pi \\nu {\\mathcal G}_0 \\vecb S \n =-\\frac{d_2}{d_1}\\nu \\vecb u^1\\,; \\qquad \\vecb F_j(\\vecb u_j) + \n \\frac{2\\pi D_0}{\\tau \\nu} S_j \\vecb e_1 = 0\\,, \\quad j=1,\\ldots, m \\,,\n\\end{equation}\nwhere the $m\\times m$ matrices ${\\mathcal E}$ and the Neumann Green's \nmatrix ${\\mathcal G}_0$, with entries $({\\mathcal G}_0)_{ij}$, are\ndefined by\n\\begin{equation}\\label{bdEJ}\n{\\mathcal E} \\equiv \\frac{1}{m} \\vecb e \\vecb e^T \\,; \\qquad\n\\left({\\mathcal G}_0\\right)_{ij} = G_{0}(\\vecb x_i;\\vecb x_j)\\equiv G_{0,ij} \\,,\n \\quad i\\neq j \\,, \\qquad\n\\left({\\mathcal G}_0\\right)_{ii} = R_{0,i} \\,.\n\\end{equation}\nHere $\\vecb e$ is the $m$-vector $\\vecb e\\equiv (1,\\ldots,1)^T$. The\nleading-order solution to (\\ref{bdsysS2m1}) when $\\nu\\ll 1$ has the\nform\n\\begin{equation}\\label{dlarge:s}\n \\vecb S = \\nu \\vecb S_0 + {\\mathcal O}(\\nu^2) \\,, \\qquad\n \\vecb u_j = \\vecb u_{j0} + {\\mathcal O}(\\nu) \\,.\n\\end{equation}\nFrom (\\ref{bdsysS2m1}) we conclude that $\\vecb S_0$ and $\\vecb u_{j0}$\nsatisfy the limiting leading-order nonlinear algebraic system\n\\begin{equation}\\label{Dlarge:1ss}\n\\left(1+\\frac{D_0}{d_1}\\right)\\vecb S_0+\\frac{2\\pi m D_0}{|\\Omega|} \n {\\mathcal E} \\vecb S_0 \n =-\\frac{d_2}{d_1} \\vecb u_0^1\\,; \\qquad \\vecb F_j(\\vecb u_{0j}) + \n \\frac{2\\pi D_0}{ \\tau} S_{0j} \\vecb e_1 = 0\\,, \\quad j=1,\\ldots, m \\,.\n\\end{equation}\nSince this leading order system does not involve the Neumann Green's\nmatrix ${\\mathcal G}_0$, we conclude that $\\vecb S_0$ is independent of the \nspatial configuration of the cells within $\\Omega$.\n\nFor the special case where the kinetics $\\vecb F_j$ is\nidentical for each cell, so that $\\vecb F_j=\\vecb F$ for\n$j=1,\\ldots,m$, we look for a solution to (\\ref{Dlarge:1ss}) with\nidentical source strengths, so that $S_{0j}$ and \n$\\vecb u_{0j}=\\vecb u_{0}$ are independent of $j$. Therefore, we write\n\\begin{equation}\n \\vecb S_0 = S_{0c} \\vecb e \\,, \\label{Dlarge:scom}\n\\end{equation}\nwhere $S_{0c}$ is the common source strength. From\n(\\ref{Dlarge:1ss}), where we use ${\\mathcal E}\\vecb e=\\vecb e$, this\nyields that $S_{0c}$ and $\\vecb u_0$ satisfy the $m+1$ dimensional\nnonlinear algebraic system\n\\begin{equation}\\label{Dlarge:ss}\n\\left(1+\\frac{D_0}{d_1} + \\frac{2\\pi m D_0}{|\\Omega|} \\right) S_{0c}\n = -\\frac{d_2}{d_1} u_0^1\\,, \\qquad \\vecb F(\\vecb u_{0}) + \n \\frac{2\\pi D_0}{ \\tau} S_{0c} \\vecb e_1 = 0\\,,\n\\end{equation}\nwhere $u_0^1$ is the first component of $\\vecb u_0$. This simple\nlimiting system will be studied in detail in the next section for\nvarious choices of the nonlinear intracellular kinetics $\\vecb F(\\vecb\nu_0)$.\n\nNext, we will simplify the GCEP, given by (\\ref{gcep:full}), when\n$D={D_0\/\\nu}\\gg 1$, and under the assumption that the reaction\nkinetics are the same in each cell. In the same way as was derived in\n(\\ref{bdgreenssD})--(\\ref{bdgexp1}), we let $D={D_0\/\\nu}\\gg 1$ and\napproximate the $\\lambda$-dependent reduced Green's function\n$G_\\lambda(\\vecb x;\\vecb x_j)$, which satisfies (\\ref{bdgneig}). For\n$\\tau={\\mathcal O}(1)$, we calculate, in place of (\\ref{bdgexp1}),\nthat the two-term expansion in terms of the Neumann Green's function\n$G_0$ is\n\\begin{equation*}\nG_\\lambda(\\vecb x;\\vecb x_j) = \\frac{D_0}{\\nu (1+\\tau\\lambda) |{\\Omega}|}\n +G_{0}(\\vecb x;\\vecb x_j)+ {\\mathcal O}(\\nu)\\,, \\qquad\nR_{\\lambda,j} = \\frac{D_0}{\\nu (1+\\tau\\lambda)|{\\Omega}|}\n +R_{0,j} + {\\mathcal O}(\\nu)\\,.\n\\end{equation*}\nTherefore, for $D={D_0\/\\nu}\\gg 1$ and $\\tau={\\mathcal O}(1)$, we have in\nterms of ${\\mathcal E}$ and the Neumann Green's matrix ${\\mathcal G}_0$ of\n(\\ref{bdEJ}), that\n\\begin{equation}\\label{dlarge:glam}\n {\\mathcal G}_\\lambda=\\frac{mD_0}{\\nu(1+\\tau\\lambda)|{\\Omega}|}\n {\\mathcal E} + {\\mathcal G}_0 + {\\mathcal O}(\\nu) \\,.\n\\end{equation}\n\nWe substitute (\\ref{dlarge:glam}) into (\\ref{gcep:mdef}), and set\n$D={D_0\/\\nu}$. In (\\ref{gcep:mdef}), we calculate to leading order in\n$\\nu$ that the matrix ${\\mathcal K}(\\lambda)$, defined in (\\ref{kdef:2}), \nreduces to \n\\begin{equation}\\label{klambda:2t}\n {\\mathcal K} \\sim \\frac{M_{11}}{\\det(\\lambda I - J)} + {\\mathcal O}(\\nu)\\,,\n\\end{equation}\nwhere $J$ is the Jacobian of $\\vecb F$ evaluated at the solution\n$\\vecb u_0$ to the limiting problem (\\ref{Dlarge:ss}), and $M_{11}$ is\nthe cofactor of $\\lambda I - J$ associated with its first row and\nfirst column. The ${\\mathcal O}(\\nu)$ correction in ${\\mathcal\n K}(\\lambda)$ arises from the higher order terms in the Jacobian\nresulting from the solution to the full system (\\ref{bdsysS2m1}). \n\nIn this way, the matrix ${\\mathcal M}$ in (\\ref{gcep:mdef}), which is\nassociated with the GCEP, reduces to leading order to\n\\begin{subequations} \\label{mdef:dlarge}\n\\begin{equation}\\label{gcep:mdef_large}\n {\\mathcal M} = a(\\lambda) I + b(\\lambda) {\\mathcal E} + {\\mathcal O}(\\nu)\n \\,,\n\\end{equation}\nwhere $a(\\lambda)$ and $b(\\lambda)$ are defined by\n\\begin{equation} \\label{mdef:dlarge_ab}\n a(\\lambda) = 1 + \\frac{D_0}{d_1} + \\frac{2\\pi d_2}{d_1 { \\tau}} \n \\frac{D_0 M_{11}}\n {\\det(\\lambda I-J)} \\,, \\qquad b(\\lambda) = \n \\frac{2\\pi m D_0}{(1+\\tau\\lambda)|\\Omega|} \\,. \n\\end{equation}\n\\end{subequations}$\\!$ We remark that the ${\\mathcal O}(\\nu)$ correction terms in\n(\\ref{gcep:mdef_large}) arises from both $2\\pi\\nu {\\mathcal G}_0$,\nwhich depends on the spatial configuration of the cells within\n$\\Omega$, and the ${\\mathcal O}(\\nu)$ term in ${\\mathcal K}$ as\nwritten in (\\ref{klambda:2t}).\n\nTherefore, when $D={D_0\/\\nu}$, it follows from (\\ref{gcep:mdef_large})\nand the criterion (\\ref{bdMdet}) of the GCEP that $\\lambda$ is a\ndiscrete eigenvalue of the linearization if and only if there exists a\nnontrivial solution $\\vecb c$ to\n\\begin{equation}\\label{dlarge:key}\n \\left( a(\\lambda) I + b(\\lambda) {\\mathcal E} \\right) \\vecb c = \n \\vecb 0\\,.\n\\end{equation}\nAny such eigenvalue with $\\mbox{Re}(\\lambda)>0$ leads to a linear instability\nof the steady-state solution when $D={\\mathcal O}(\\nu^{-1})$.\n\nWe now derive explicit stability criteria from (\\ref{dlarge:key}) by\nusing the key properties that ${\\mathcal E}\\vecb e=\\vecb e$ and\n${\\mathcal E} \\vecb q_j=0$ for $j=2,\\ldots,m$, where $\\vecb q_j$ for\n$j=2,\\ldots,m$ are an orthogonal basis of the $m-1$ dimensional\nperpendicular subspace to $\\vecb e$, i.e $\\vecb q_j^T\\vecb e=0$.\nWe obtain that $\\lambda$ is a discrete eigenvalue for \nthe {\\em synchronous mode}, corresponding to $\\vecb c=\\vecb e$, whenever\n$\\lambda$ satisfies\n\\begin{equation}\\label{dlarge:sync}\n a(\\lambda)+b(\\lambda) \\equiv \n 1 + \\frac{D_0}{d_1} + \\frac{2\\pi d_2}{d_1 { \\tau}} \\frac{D_0 M_{11}}\n {\\det(\\lambda I-J)} + \\frac{2\\pi m D_0}{(1+\\tau\\lambda)|\\Omega|} =0 \\,.\n\\end{equation}\nThis expression can be conveniently written as\n\\begin{equation}\\label{dlarge:sync_new}\n \\frac{M_{11}}{\\det(\\lambda I-J)} = -\\frac{{ \\tau}}{2\\pi d_2} \\left( \n \\frac{\\kappa_1 \\tau \\lambda + \\kappa_2}{\\tau\\lambda + 1} \\right) \\,, \\qquad\n \\mbox{where} \\qquad \n \\kappa_1 \\equiv \\frac{d_1}{D_0} + 1 \\,, \\quad \\kappa_2 \\equiv \\kappa_1 + \n \\frac{2m \\pi d_1}{|\\Omega|} \\,. \n\\end{equation}\nIn contrast, $\\lambda$ is a discrete eigenvalue for \nthe {\\em asynchronous or competition modes}, corresponding to \n$\\vecb c=\\vecb q_j$ for $j=2,\\ldots,m$, whenever $\\lambda$ satisfies\n$a(\\lambda)=0$. This yields for any $m\\geq 2$ that\n\\begin{equation}\\label{dlarge:async}\n \\frac{M_{11}}{\\det(\\lambda I-J)} = - \\frac{{ \\tau}}{2\\pi\n d_2} \\left( \\frac{d_1}{D_0} + 1 \\right) \\,.\n\\end{equation}\n\nAny discrete eigenvalue for either of the two modes that\nsatisfies $\\mbox{Re}(\\lambda)>0$ leads to an instability. If all such\neigenvalues satisfy $\\mbox{Re}(\\lambda)<0$, then the steady-state\nsolution for the regime $D={D_0\/\\nu}$ is linearly stable on an\n${\\mathcal O}(1)$ time-scale.\n\n\\begin{remark}\\label{remark} The spectral problems (\\ref{dlarge:sync_new}) and\n(\\ref{dlarge:async}) have a remarkably similar form to the spectral\n problem characterizing the linear stability of localized spot\n solutions to the Gierer-Meinhardt (GM) RD system\n\\begin{equation}\\label{gm}\n a_t=\\epsilon^2 \\Delta a - a + {a^2\/h} \\,, \\qquad \\tau h_t=D \\Delta h-\nh + \\epsilon^{-2} a^2 \\,,\n\\end{equation}\nposed in a bounded two-dimensional domain $\\Omega$ with $\\partial_n\na=\\partial_n h=0$ on $\\partial\\Omega$. For $\\epsilon\\to 0$, and for\n$D={D_0\/\\nu}$ with $\\nu={-1\/\\log\\epsilon}\\ll 1$, the linear stability\nof an $m$-spot solution, with eigenvalue parameter $\\lambda$,\nis characterized by the roots of a transcendental equation of the\nform (see \\cite{wei1})\n\\begin{equation}\\label{gm:stab}\n \\frac{a+b \\tau\\lambda}{c+d \\tau \\lambda} = {\\mathcal F}(\\lambda) \\,,\n \\qquad {\\mathcal F}(\\lambda)\\equiv \\frac{\\int_{0}^{\\infty} w\\left(L_0-\\lambda\n \\right)^{-1} w^2 \\rho \\, d\\rho}{\\int_{0}^{\\infty} w^2\\rho\\, d\\rho} \\,,\n\\end{equation}\nwhere $w(\\rho)>0$ is the radially symmetric ground-state solution of\n$\\Delta w-w+w^2=0$ with $w\\to 0$ as $\\rho\\to\\infty$, and $L_0$ is the\nlocal operator $L_0\\Phi\\equiv \\Delta \\Phi - \\Phi + 2w\\Phi$, where\n$L_0^{-1}$ is restricted to radially symmetric functions. In direct\nanalogy to our cell-bulk problem, it was shown in \\cite{wei1} that\nthere can be either synchronous or asynchronous instabilities of the\namplitudes of the localized spots. For synchronous instabilities, the\ncoefficients in (\\ref{gm:stab}) are $a=b=1$, $c=2$, and\n$d={2\/(1+\\mu)}$, where $\\mu\\equiv {2\\pi m D_0\/|\\Omega|}$, while for\nasynchronous instabilities we have $a={(1+\\mu)\/2}$, and\n$b=c=d=0$. Since $L_0$ has a unique eigenpair in $H^{1}$ with positive\neigenvalue $\\sigma_0>0$, the function ${\\mathcal F}(\\lambda)$ in\n(\\ref{gm:stab}) is analytic in $\\mbox{Re}(\\lambda)>0$, with the\nexception of a simple pole at $\\lambda=\\sigma_0>0$. In this way,\n(\\ref{gm:stab}) is remarkably similar to our cell-bulk spectral\nproblems (\\ref{dlarge:sync_new}) and (\\ref{dlarge:async}) when there\nis only a single intracellular species that is self-activating in the\nsense that $J\\equiv F_{u}(u_e)>0$ so that ${M_{11}\/\\det(\\lambda\n I-J)}={1\/(\\lambda-F_u(u_e))}$. Spectral problems similar to\n(\\ref{gm:stab}) also occur for the Schnakenberg RD\nsystem \\cite{wei2}, the Brusselator \\cite{rozada}, and the Gray-Scott\nmodel \\cite{chen}. For the GM model, it can be shown that there is a\nunique Hopf bifurcation value of $\\tau$ for the synchronous mode\nwhenever $\\mu>1$. However, for our cell-bulk model, below in \\S\n\\ref{one:ode} we prove that Hopf bifurcations are impossible for the\nsynchronous mode whenever there is only one intracellular species.\n\\end{remark}\n\n\n\\setcounter{equation}{0}\n\\setcounter{section}{3}\n\\section{Examples of the Theory: Finite Domain With $D={\\mathcal O}(\\nu^{-1})$}\n\\label{sec:examples}\n\nIn this section we will study the leading-order steady-state problem\n(\\ref{Dlarge:ss}), and its associated spectral problem\n(\\ref{dlarge:sync_new}) and (\\ref{dlarge:async}), for various special\ncases and choices of the reaction kinetics $\\vecb F$. We investigate\nthe stability properties of the steady-state (\\ref{Dlarge:ss}) as\nparameters are varying, and in particular, find conditions for Hopf\nbifurcations to occur.\n\n\\subsection{Example 1: $m$ Cells; One Local Component}\\label{one:ode}\nTo illustrate our theory, we first consider a system such that the\nlocal dynamics inside each cell consists of only a single component\nwith arbitrary scalar kinetics $F$. For this case, the steady-state\nproblem for $\\vecb u_0=u_0$ and $S$, given by (\\ref{Dlarge:ss}),\nreduces to two algebraic equations. We will study the stability\nproblem for both the synchronous and asynchronous modes. We show that\nfor any $F(u)$, the steady-state can never be destablilized by a Hopf\nbifurcation.\n\nFor the one-component case, we calculate $M_{11}=1$ and $\\det(\\lambda\nI-J)=\\lambda-F_u^e$, where $F_u^e$ is defined as the derivative of\n$F(u)$ evaluated at the steady-state $u_0$. From (\\ref{dlarge:sync_new}),\nthe spectral problem for the synchronous mode reduces to \n\\begin{subequations} \\label{d0:one}\n\\begin{equation}\\label{mcell1ode:sync}\n\\lambda^2-\\lambda p_1 + p_2 =0 \\,, \\qquad p_1 \\equiv\nF_u^e-\\frac{\\gamma}{ \\tau} -\\frac{\\zeta}{\\tau} \\,, \\qquad p_2 \\equiv\n\\frac{1}{\\tau}\\left(\\frac{\\gamma}{ \\tau}-\\zeta F_u^e\\right) \\,,\n\\end{equation}\nwhere\n\\begin{equation}\\label{alphabeta}\n\\gamma\\equiv\\frac{2\\pi d_2 D_0}{d_1+D_0}>0\\,, \\qquad\n\\zeta\\equiv1+\\frac{2\\pi md_1D_0}{|\\Omega| (d_1+D_0)}>1\\,.\n\\end{equation}\n\\end{subequations}$\\!$ The following result characterizes the stability properties for\nthe synchronous mode:\n\n\\begin{result}\\label{d0:theorem_1}\nThere can be no Hopf bifurcations associated with the synchronous mode.\nMoreover, suppose that\n\\begin{equation}\\label{d0:stab_one}\n F_{u}^e < \\frac{\\gamma}{\\zeta \\tau} = \\frac{2\\pi\n d_2}{ \\tau} \\left[ 1+ \\frac{d_1}{D_0} + \\frac{2\\pi m\n d_1}{|\\Omega|} \\right]^{-1} \\,.\n\\end{equation}\nThen, we have $\\mbox{Re}(\\lambda)<0$, and so the steady-state is\nlinearly stable to synchronous perturbations. If $F_{u}^{e} >\n{\\gamma\/\\left(\\zeta \\tau\\right)}$, the linearization has\nexactly one positive eigenvalue.\n\\end{result}\n\n\\begin{proof}\nFor a Hopf bifurcation to occur we need $p_1=0$ and $p_2>0$. Upon\nsetting $p_1=0$, we get $F_u^e={\\left(\\gamma+ \\zeta\\right)\/\\tau} >0$.\nUpon substituting this expression into the formula for\n$p_2$ in (\\ref{mcell1ode:sync}) we get\n\\begin{equation}\np_2 = \\frac{1}{ \\tau}\\left( \\frac{\\gamma}{ \\tau}\n -\\zeta F_u^e\\right)=\\frac{1}{{\\tau^2}} \\left( \\gamma(1-\\zeta)-\n\\zeta^2\\right)<0\\,,\n\\end{equation}\nsince $\\gamma>0$ and $\\zeta>1$. Therefore, there can be no Hopf\nbifurcation for the synchronous mode.\n\nNext, to establish the stability threshold, we note that the\nsteady-state solution is stable to synchronous perturbations if and\nonly if $p_1<0$ and $p_2>0$. From (\\ref{mcell1ode:sync}), we have that\n$p_1<0$ and $p_2>0$ when\n\\begin{equation}\\label{equal}\n \\tau F_u^e < \\zeta + \\gamma \\,, \\qquad \\tau F_u^e < {\\gamma\/\\zeta}\\,,\n\\end{equation}\nrespectively, which implies that we must have $\\tau F_u^e\n<\\min(\\zeta+\\gamma,{\\gamma\/\\zeta})$. Since $\\zeta>1$, the two\ninequalities in (\\ref{equal}) hold simultaneously only when $\\tau\nF_u^e <{\\gamma\/\\zeta}$. This yields that\n$\\mbox{Re}(\\lambda)<0$ when (\\ref{d0:stab_one}) holds. Finally, when\n$F_{u}^{e} >{\\gamma\/\\left(\\zeta \\tau\\right)}$, then $p_2<0$,\nand so there is a unique positive eigenvalue. \\end{proof}\n\nThis result shows that the effect of cell-bulk coupling is that the\nsteady-state of the coupled system can be linearly stable even when\nthe reaction kinetics is self-activating in the sense that\n$F_{u}^{e}>0$. We observe that the stability threshold\n${\\gamma\/\\zeta}$ is a monotone increasing function of $D_0$, with\n${\\gamma\/\\zeta}\\to 0$ as $D_0\\to 0$ and ${\\gamma\/\\zeta}$ tending to a\nlimiting value as $D_0\\to \\infty$. This shows that as $D_0$ is\ndecreased, corresponding to when the cells are effectively more\nisolated from each other, there is a smaller range of $F_u^e>0$ where\nstability can still be achieved.\n\nNext, we will consider the spectral problem for the asynchronous mode.\nFrom (\\ref{dlarge:async}), we get\n\\begin{equation}\n\\frac{1}{\\lambda-F_u^e}=-\\frac{{\\tau}}{\\gamma}\\,,\n\\end{equation}\nwhere $\\gamma$ is defined in (\\ref{alphabeta}). Therefore,\n$\\lambda=F_u^e-{\\gamma\/{\\tau}}$, and so $\\lambda$ is real and\nno Hopf bifurcation can occur. This asynchronous mode is stable if\n$F_u^e<{\\gamma\/ \\tau}$. Since $\\zeta>1$, we observe, upon\ncomparing this threshold with that for the synchronous mode in\n(\\ref{d0:stab_one}), that the stability criterion for the synchronous\nmode is the more restrictive of the two stability thresholds.\n\nIn summary, we conclude that a Hopf bifurcation is impossible for\n(\\ref{mainbd}) in the parameter regime $D={D_0\/\\nu}$ when there is only\none dynamically active species inside each of $m$ small cells.\nMoreover, if $F_{u}^e < {\\gamma\/(\\zeta \\tau)}$, where\n$\\gamma$ and $\\zeta$ are defined in (\\ref{alphabeta}), then the\nsteady-state solution is linearly stable to both the synchronous and\nasynchronous modes.\n\n\\subsection{Example 2: $m$ Cells; Two Local Components}\n\nNext, we assume that there are two dynamically active local species\ninside each of $m$ distinct cells. For ease of notation, we write the\nintracellular variable as $\\vecb u=(v,w)^T$ and the local kinetics as\n$\\vecb F(v,w)=(F(v,w),G(v,w))^T$. In this way, the steady-state\nproblem (\\ref{Dlarge:ss}) becomes\n\\begin{equation}\\label{mcell2ode:ss}\n\\left(1+\\frac{D_0}{d_1} + \\frac{2\\pi m D_0}{|\\Omega|} \\right) S_{0c} =\n-\\frac{d_2}{d_1} v_e\\,, \\qquad F(v_e,w_e)+\\frac{2\\pi D_0}{\n \\tau} S_{0c}= 0\\,, \\qquad G(v_e,w_e)=0\\,.\n\\end{equation}\nGiven specific forms for $F$ and $G$, we can solve the steady-state\nproblem (\\ref{mcell2ode:ss}) either analytically or numerically.\n\nTo analyze the stability problem, we first calculate the cofactor\n$M_{11}$ as $M_{11}=\\lambda-G_w^e$ and $\\det(\\lambda\nI-J)=\\lambda^2-\\mbox{tr}(J)\\lambda+\\det(J)$, where $\\mbox{tr}(J)$ and\n$\\det(J)$ are the trace and determinant of the Jacobian of $\\vecb F$,\ngiven by\n\\begin{equation}\n \\mbox{tr}(J)=F_v^e+G_w^e\\,,\\qquad \\det(J)=F_v^eG_w^e-F_w^eG_v^e\\,.\n\\end{equation}\nHere $F_i^e$, $G_i^e$ are partial derivatives of $F$, $G$ with respect\nto $i$, with $i\\in(v,w)$, evaluated at the solution to\n(\\ref{mcell2ode:ss}).\n\nNext, we analyze the stability of the steady-state solution with respect to\neither synchronous or asynchronous perturbations. For the synchronous mode, \nwe obtain, after some algebra, that (\\ref{dlarge:sync_new}) can be\nreduced to the study of the cubic\n\\begin{subequations} \\label{d0:cubic}\n\\begin{equation}\\label{mcell2ode:H3}\n\\mathcal{H}(\\lambda)\\equiv\\lambda^3+\\lambda^2 p_1+\\lambda p_2+p_3=0\\,.\n\\end{equation}\nwhere $p_1$, $p_2$, and $p_3$, are defined in terms of\n$\\gamma$ and $\\zeta$, as given in (\\ref{alphabeta}), by\n\\begin{equation}\\label{mcell2ode:p123}\n\\begin{aligned}\np_1\\equiv \\frac{\\gamma}{ \\tau}\n+\\frac{\\zeta}{\\tau}-\\mbox{tr}(J)\\,,\\qquad p_2\\equiv\n\\det(J)-\\frac{\\gamma}{ \\tau}\nG_w^e+\\frac{1}{\\tau}(\\frac{\\gamma}{ \\tau}-\\zeta\n\\mbox{tr}(J))\\,,\\qquad p_3\\equiv \\frac{1}{\\tau}(\\zeta \\det(J)-\n \\frac{\\gamma}{ \\tau} G_w^e)\\,.\n\\end{aligned}\n\\end{equation}\n\\end{subequations}$\\!$\n\nTo determine whether there is any eigenvalue in\n$\\mbox{Re}(\\lambda)>0$, and to detect any Hopf bifurcation boundary in\nparameter space, we use the well-known \\textit{Routh-Hurwitz\n criterion} for a cubic function. It is well known that all three\nroots to the cubic satisfy $\\mbox{Re}(\\lambda)<0$ if and only if the\nfollowing three conditions on the coefficients hold:\n\\begin{equation}\np_1>0\\,, \\qquad p_3>0\\,,\\qquad p_1p_2>p_3\\,.\n\\end{equation}\nTo find the Hopf bifurcation boundary, we need only consider a special\ncubic equation which has roots $\\lambda_1=a<0$ and $\\lambda_{2,3}=\\pm\ni\\omega$, for which\n$(\\lambda-a)(\\lambda-i\\omega)(\\lambda+i\\omega)=\\lambda^3-\na\\lambda^2+\\omega^2\\lambda-a\\omega^2=0\\,.$ Comparing this expression\nwith (\\ref{mcell2ode:H3}), and using the \\textit{Routh-Hurwitz\n criterion}, we conclude that on any Hopf bifurcation boundary the\nparameters must satisfy\n\\begin{equation}\np_1>0\\,,\\qquad p_3>0\\,, \\qquad p_1p_2=p_3\\,. \\label{mcell2ode:syn}\n\\end{equation} \nWe will return to this criterion below when we study two specific\nmodels for the local kinetics $(F,G)$.\n\nNext, we consider the spectral problem for the asynchronous mode. Upon\nsubstituting the expressions of $M_{11}$ and $\\det(\\lambda I-J)$ into\n(\\ref{dlarge:async}) and reorganizing the resulting expression,\n(\\ref{dlarge:async}) becomes the quadratic equation \n\\begin{equation}\n\\lambda^2-\\lambda q_1 + q_2 = 0 \\,, \\qquad \\mbox{where} \\qquad q_1\n\\equiv \\mbox{tr}(J)- \\frac{\\gamma}{ \\tau} \\,, \\qquad q_2 \\equiv\n\\det(J)- \\frac{\\gamma}{ \\tau} G_w^e \\,.\n\\end{equation}\nFor a Hopf bifurcation to occur, we require that $q_1=0$ and $q_2>0$, which\nyields that\n\\begin{equation}\\label{mcell2ode:asyn_all}\n\\frac{\\gamma}{ \\tau} =\\mbox{tr}(J)=F_v^e+G_w^e\\,, \\qquad \\mbox{provided that}\n \\qquad \\det(J)-\\frac{\\gamma}{ \\tau} G_w^e=-G_v^eF_w^e-(G_w^e)^2>0\\,.\n\\end{equation}\nFinally, we conclude that $\\mbox{Re}(\\lambda)<0$ for the\nasynchronous modes if and only if\n\\begin{equation}\n \\mbox{tr}(J)<{\\gamma\/ \\tau} \\,, \\qquad \\mbox{and} \\qquad \n \\det(J) -\\frac{\\gamma}{ \\tau} G_w^e>0\n \\,.\n\\end{equation}\n\nTo write the stability problem for the asynchronous mode in terms of\n$D_0$, we use (\\ref{alphabeta}) for $\\gamma$ in terms of $D_0$ to\nobtain from the conditions (\\ref{mcell2ode:asyn_all}) that the Hopf\nbifurcation threshold satisfies the transcendental equation\n\\begin{equation}\\label{mcell:d0_async}\nD_0=\\frac{\\tau d_1 \\mbox{tr}(J)}{2\\pi d_2-\\tau \\mbox{tr}(J)}\\,, \\qquad\n \\mbox{provided that} \\qquad \n D_0\\left( \\frac{2\\pi d_2}{\\tau} G_w^e-\\det (J)\\right)0$, $\\mu>0$, and $\\epsilon_0>0$ by\n\\begin{equation}\\label{2dsel:fg}\nF(v,w)=\\alpha w+wv^2-v\\,, \\qquad G(v,w) =\\epsilon_0\\left\n(\\mu-(\\alpha w+wv^2)\\right) \\,.\n\\end{equation}\nFirst, we determine the approximate steady-state solution $v_e$ and\n$w_e$ by substituting (\\ref{2dsel:fg}) into (\\ref{mcell2ode:ss}). This\nyields that\n\\begin{equation}\\label{2dsel:ss_sol}\n v_e=\\frac{\\mu}{\\left[1+ {2\\pi D_0\\beta\/{ \\tau}}\\right]}\\,, \\qquad\n w_e=\\frac{\\mu}{\\alpha+v_e^2}\\,, \\qquad S_{0c}= -\\beta v_e\\,, \\qquad\n \\mbox{where} \\qquad \\beta \\equiv \\frac{d_2}{d_1+D_0+2\\pi md_1D_0\/|\\Omega|}\\,.\n\\end{equation}\nAs needed below, the partial derivatives of $F$ and $G$ evaluated at\nthe steady-state solution are\n\\begin{subequations} \\label{2dsel:jac}\n\\begin{equation}\\label{2dsel:fgp}\nF_v^e = 2v_ew_e-1\\,, \\qquad F_w^e=\\alpha+v_e^2\\,, \\qquad\nG_v^e=-2\\epsilon_0 v_ew_e\\,, \\qquad G_w^e=-\\epsilon_0(\\alpha+v_e^2)\\,,\n\\end{equation}\nwhich yields\n\\begin{equation}\n \\mbox{det}(J) = \\epsilon_0\\left(\\alpha+ v_e^2\\right)=-G_w^e >0 \\,, \\qquad\n \\mbox{tr}(J)= 2 v_e w_e -1 -\\epsilon_0 \\left(\\alpha + v_e^2\\right) \\,.\n\\end{equation}\n\\end{subequations}$\\!$\n\nTo study possible synchronous oscillations of the $m$ cells, we\ncompute the Hopf bifurcation boundaries as given in\n(\\ref{mcell2ode:syn}), where we use (\\ref{2dsel:jac}). For the\nparameter set $\\tau=1$, $D_0=1$, $|\\Omega|=10$, $\\mu=2$, $\\alpha=0.9$, and\n$\\epsilon_0=0.15$, we obtain the Hopf bifurcation boundary in the\n$d_1$ versus $d_2$ parameter plane as shown by the solid curves in\nFig.~\\ref{d1d2sel} for $m=1,2,3$.\n\nNext, to obtain instability thresholds corresponding to the asynchronous\nmode, we substitute (\\ref{2dsel:jac}) into (\\ref{mcell2ode:asyn_all})\nto obtain that the Hopf bifurcation boundary is given by\n\\begin{equation}\n\\gamma=\\tau \\mbox{tr}(J) \\equiv \\tau\n \\left[2v_ew_e-1-\\epsilon(\\alpha+v_e^2) \\right]\\,,\n\\label{2dsel:asy_1}\n\\end{equation}\nprovided that $\\mbox{det}(J)-\\mbox{tr}(J) G_{w}^e>0$. This latter\ncondition can be written, using (\\ref{2dsel:jac}), as\n$\\epsilon_0(\\alpha+v_e^2)\\left(1+\\mbox{tr}(J)\\right)>0$, and so is\nsatisfied provided that $\\mbox{tr}(J)>-1$. Since $\\gamma>0$ from\n(\\ref{alphabeta}), (\\ref{2dsel:asy_1}) implies that we must have\n$\\mbox{tr}(J)>0$, which guarantees that $\\mbox{det}(J)-\\mbox{tr}(J)\nG_{w}^e>0$ always holds at a Hopf bifurcation. In this way, and by\nsubstituting (\\ref{2dsel:ss_sol}) for $w_e$ into (\\ref{2dsel:asy_1}),\nwe obtain that the asynchronous mode has a pure imaginary pair of\ncomplex conjugate eigenvalues when\n\\begin{equation}\\label{2dsel:hopf_async}\n\\gamma= \\tau \\left[\\frac{2 v_e \\mu}{\\alpha + v_e^2} -1 -\n \\epsilon_0 \\left(\\alpha + v_e^2 \\right) \\right] \\,, \\qquad \\mbox{where}\n \\qquad v_e = \\frac{\\mu}{\\left[1+ {2\\pi D_0\\beta\/ \\tau}\\right]}\\,.\n\\end{equation}\nHere $\\gamma$ and $\\beta$, depending on $d_1$, $d_2$, $m$,\n$|\\Omega|$, and $D_0$, are defined in (\\ref{alphabeta}) and\n(\\ref{2dsel:ss_sol}), respectively. By using these expressions for\n$\\gamma$ and $\\beta$, we can readily determine a parametric form for\nthe Hopf bifurcation boundary in the $d_1$ versus $d_2$ plane as the\nsolution to a linear algebraic system for $d_1$ and $d_2$ in terms of\nthe parameter $v_e$ with $00$ and $d_2>0$, we obtain a parametric form for\nthe Hopf bifurcation boundary for the asynchronous mode in the \n$d_1$ versus $d_2$ parameter plane. For $m=2$ and $m=3$, these are\nthe dashed curves shown in Fig.~\\ref{d1d2sel}.\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{sel_hopfbd.eps}\n\\caption{Hopf bifurcation boundaries for the synchronous (solid curve)\n and asynchronous (dashed curve) modes for the Sel'kov model for\n different numbers $m$ of cells in the $d_1$ versus $d_2$ parameter\n plane. The synchronous mode for $m=1$ is unstable between the two\n black lines. For $m=2$ and $m=3$ the synchronous mode is\n unstable in the horseshoe-shaped region bounded by the blue and red\n solid curves, respectively. Inside the dotted regions for $m=2$ and\n $m=3$ the asynchronous mode is unstable. For the asynchronous mode,\n the boundary of these regions is given parametrically by\n (\\ref{2dcell:a_par}). The parameters used are $\\mu=2$, $\\alpha=0.9$,\n $\\epsilon_0=0.15$, $\\tau=1$, $D_0=1$, and $|\\Omega|=10$.}\n\\label{d1d2sel}\n\\end{center}\n\\end{figure}\n\nWe now discuss qualitative aspects of the Hopf bifurcation boundaries\nfor both synchronous and asynchronous modes for various choices of $m$\nas seen in Fig.~\\ref{d1d2sel}. For $m=1$, we need only consider the\nsynchronous instability. The Hopf bifurcation boundary is given by the\ntwo black lines, and the region with unstable oscillatory dynamics is\nlocated between these two lines. For $m=2$, inside the region bounded\nby the blue solid curve, the synchronous mode is unstable and under\nthe blue dashed curve, the asynchronous mode is unstable. Similar\nnotation applies to the case with $m=3$, where the Hopf bifurcation\nboundaries for synchronous\/asynchronous mode are denoted by red\nsolid\/dashed curves.\n\nOne key qualitative feature we can observe from Fig.~\\ref{d1d2sel}, for\nthe parameter set used, is that the oscillatory region for a larger\nvalue of $m$ lies completely within the unstable region for smaller\n$m$ for both the synchronous and asynchronous modes. This suggests\nthat if a coupled system with $m_1$ cells is unstable to synchronous\nperturbations, then a system with $m_20$, $q>0$, and $z>1$. \n\nUpon substituting (\\ref{2dsf:fg}) into (\\ref{mcell2ode:ss}) we\ncalculate that the steady-state solution $w_e>0$ is given by the unique\nreal positive root of the cubic ${\\mathcal C}(w)=0$ given by\n\\begin{equation}\n {\\mathcal C}(w)\\equiv qw^3-6qw^2+w\\left(12q-1+\\Lambda\\right)-(8q+4)\\,,\\qquad\n \\mbox{where} \\qquad\n\\Lambda\\equiv\\frac{\\epsilon_0 z}{\\left[\\epsilon_0 + {2\\pi\n D_0\\beta\/ \\tau} \\right]}\\,, \\label{2dfitz:cubic}\n\\end{equation}\nwhere $\\beta$ is defined in (\\ref{2dsel:ss_sol}). The uniqueness of the\npositive root of this cubic for any $\\Lambda>0$ was proved previously\nin \\S 2 of \\cite{gou_ward}. In terms of the solution $w_e>0$ to the\ncubic equation, we calculate $v_e=\\Lambda w_e$ and $S_{0c}=-\\beta \\Lambda w_e$.\n\nAs needed below, we first calculate the partial derivatives of $F$ and $G$\nat the steady-state solution as\n\\begin{subequations} \\label{2dsf:jac}\n\\begin{equation}\\label{2dsf:fgp}\nF_v^e = -\\epsilon_0\\,, \\qquad F_w^e=\\epsilon_0 z\\,, \\qquad\nG_v^e=-1\\,, \\qquad G_w^e=1 - 3q(w_e-2)^2 \\,,\n\\end{equation}\nwhich yields\n\\begin{equation}\n \\mbox{det}(J) = \\epsilon_0\\left[ z-1 + 3q(w_e-2)^2\\right]>0 \\,, \\qquad\n \\mbox{tr}(J)= -\\epsilon_0 + 1 - 3q(w_e-2)^2 \\,.\n\\end{equation}\n\\end{subequations}$\\!$\n\nTo determine conditions for which the synchronous mode has a Hopf\nbifurcation we first substitute (\\ref{2dsf:fgp}) into\n(\\ref{mcell2ode:p123}). The Hopf bifurcation boundary is then found by\nnumerically computing the roots of (\\ref{mcell2ode:syn}).\nSimilarly, to study instabilities associated with the asynchronous\noscillatory mode we substitute (\\ref{2dsf:fgp}) into\n(\\ref{mcell2ode:asyn_all}) to obtain that \n\\begin{subequations} \\label{2dfitz:async}\n\\begin{equation}\n\\gamma = \\tau \\left[-\\epsilon_0 +1-3q(w_e-2)^2\\right]\\,,\n\\end{equation}\nwhich yields the Hopf bifurcation boundary, provided that\n\\begin{equation}\\label{2dsf:asy1}\n\\det(J)-\\frac{\\gamma}{ \\tau} G_w^e=-G_v^eF_w^e-(G_w^e)^2 =\n-G_v^eF_w^e-(G_w^e)^2=\\epsilon_0 z- \\left[1-3q(w_e-2)^2\\right]^2>0\\,,\n\\end{equation}\n\\end{subequations}$\\!$ where $w_e$ is the positive root of the cubic\n(\\ref{2dfitz:cubic}). As was done for the Sel'kov model in \\S\n\\ref{sec:ex:selkov}, the Hopf bifurcation boundary for the\nasynchronous mode in the $d_1$ versus $d_2$ parameter plane can be\nparametrized as in (\\ref{2dcell:a_par_1}) where $a_{11}$, $a_{12}$,\n$a_{22}$, and $a_{21}$, are now defined in terms of the parameter\n$w_e>0$ by \\begin{subequations}\n\\begin{equation}\n a_{11}\\equiv 1 + \\frac{2\\pi m D_0}{|\\Omega|} \\,, \\qquad\n a_{12}\\equiv -\\frac{1}{\\beta(w_e)} \\,, \\qquad a_{21}\\equiv 1 \\,, \\qquad\n a_{22}\\equiv -\\frac{2\\pi D_0}{\\gamma(w_e)} \\,,\n\\end{equation}\nwhere\n\\begin{equation}\n \\beta(w_e) \\equiv \\frac{\\epsilon_0\\tau}{2\\pi D_0}\n \\left( \\frac{z}{\\Lambda(w_e)} -1 \\right) \\,, \\quad \\mbox{with}\n \\quad \\Lambda(w_e)\\equiv -\\frac{q(w_e-2)^3}{w_e} + 1 +\n \\frac{4}{w_e} \\,; \\quad \\gamma(w_e) \\equiv \\tau\n \\left[-\\epsilon_0 + 1 - 3q(w_e-2)^2 \\right]\\,.\n\\end{equation}\n\\end{subequations}$\\!$\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.75\\textwidth,height=5.0cm]{d1d2_slowfast.eps}\n\\caption{Hopf bifurcation boundaries for the synchronous (solid curve)\n and asynchronous (dashed curve) modes for the FN system \n (\\ref{2dsf:fg}) with various number $m$ of cells in the $d_1$ versus\n $d_2$ parameter plane. Between the solid lines the synchronous mode\n is unstable, whereas between the dashed lines the asynchronous mode\n is unstable. Notice that the region of instability for the asynchronous\n mode is contained within the instability region for the synchronous mode.\n Parameters used are $z=3.5$, $q=5$, $\\epsilon_0=0.5$, $\\tau=1$, $D_0=1$, \n and $|\\Omega|=10$.}\n\\label{d1d2slowfast}\n\\end{center}\n\\end{figure}\n\nBy varying $w_e>0$ and retaining only the portion of the\ncurve for which $d_1>0$ and $d_2>0$, and ensuring that the constraint\n(\\ref{2dsf:asy1}) holds, we obtain a parametric form for\nthe Hopf bifurcation boundary for the asynchronous mode in the \n$d_1$ versus $d_2$ parameter plane. For $m=2$ and $m=3$, these are\nthe dashed curves shown in Fig.~\\ref{d1d2slowfast}.\n\nIn this way, in Fig.~\\ref{d1d2slowfast} we plot the Hopf bifurcation\nboundaries for the synchronous mode (solid curves) and asynchronous\nmode (dashed curves) for various values of $m$ for the parameter set\n$z=3.5$, $q=5$, $\\epsilon_0=0.5$, $\\tau=1$, $D_0=1$, and\n$|\\Omega|=10$. As compared to Fig.~\\ref{d1d2sel}, we notice that the\nunstable regions for both modes are not only shrinking but also the\nboundaries shift as the number of cells increases. This feature\ndoes not appear in the previous Sel'kov model.\n\nNext, in Fig.~\\ref{tauD0slowfast} we show the region where oscillatory\ninstabilities occur for either synchronous or asynchronous modes for\n$m=1,2,3$ in the $\\tau$ versus $D_0$ plane. These Hopf bifurcation\nboundaries are computed by finding roots of either\n(\\ref{mcell2ode:syn}) or (\\ref{2dfitz:async}) for various values of\n$D_0$. The other parameter values are the same as those used for\nFig.~\\ref{d1d2slowfast} except $d_1=10$ and $d_2=0.2$. Inside the\nregion bounded by the solid curves the synchronous mode is unstable,\nwhile inside the region bounded by the dashed curves, the asynchronous\nmode is unstable. Similar to that shown in Fig.~\\ref{d1d2slowfast},\nthe regions of instability are shrinking, at the same time as the Hopf\nbifurcation boundaries shift, as $m$ increases. For these parameter\nvalues of $d_1$ and $d_2$, the Hopf bifurcation for the synchronous\nmodel still occurs for large values of $D_0$.\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.8\\textwidth,height=5.0cm]{tauD_slowfast.eps}\n\\caption{Hopf bifurcation boundaries for the synchronous (solid curves)\n and asynchronous (dashed curves) modes for the FN system \n (\\ref{2dsf:fg}) with various number $m$ of cells in the $\\tau$ versus\n $D_0$ parameter plane. Between the solid lines the synchronous mode\n is unstable, whereas between the dashed lines the asynchronous mode\n is unstable. Parameters used are $z=3.5$, $q=5$, $\\epsilon_0=0.5$, $d_1=10$,\n $d_2=0.2$, \n and $|\\Omega|=10$.}\n\\label{tauD0slowfast}\n\\end{center}\n\\end{figure}\n\n\n\n\\setcounter{equation}{0}\n\\setcounter{section}{4}\n\\section{Finite Domain: Reduction to ODEs for $D\\gg {\\mathcal O}(\\nu^{-1})$}\n\\label{sec:odes}\n\nIn this section we assume that there is one small dynamically active\ncircular cell $\\Omega_\\epsilon$ in the bounded domain $\\Omega$ under\nthe assumption that $D\\gg {\\mathcal O}(\\nu^{-1})$, with\n$\\nu={-1\/\\log\\epsilon}$. In this limit, in which the bulk region is\neffectively well-mixed, we now show that we can reduce the dynamics for the\ncoupled cell-bulk model to a system of nonlinear ODEs.\n\nFor the case of one cell, the basic model is formulated as\n\\begin{subequations}\\label{d:mainbd}\n\\begin{equation}\\label{d:mainU}\n\\tau U_t = D\\Delta U -U\\,, \\qquad \\vecb x\\in \\Omega\\backslash\n\\Omega_{\\epsilon}\\,; \\qquad \\partial_nU = 0\\,,\\qquad \\vecb x\\in\\partial \\Omega\\,;\n \\qquad\n\\epsilon D \\partial_{n} U =d_1U-d_2u^1\\,, \\;\\;\\quad \\vecb x\\in\\partial \n\\Omega_{\\epsilon}\\,,\n\\end{equation}\nwhich is coupled to the intracellular cell dynamics, with $\\vecb u =\n(u_1,\\ldots,u_n)^T$ and reaction kinetics $\\vecb F(\\vecb u)$, by\n\\begin{equation}\\label{d:mainuj}\n\\frac{d \\vecb u}{dt} = \\vecb F(\\vecb u)+ \\frac{\\vecb e_1}{\\epsilon \\tau}\n \\int_{\\partial \\Omega_{\\epsilon}}( d_1 U -d_2 u_1)\\, ds \\,, \\qquad\n\\mbox{where} \\qquad \\vecb e_1\\equiv (1,0,\\dots,0)^T \\,.\n\\end{equation}\n\\end{subequations}$\\!$ \n\nWe will assume that $D\\gg 1$, and so in the bulk region we expand\n\\begin{equation}\\label{d:uout}\n U=U_0+\\frac{1}{D}U_1+ \\cdots \\,.\n\\end{equation}\nUpon substituting this expansion into (\\ref{d:mainU}), and noting that\n$\\Omega_\\epsilon\\to \\vecb x_0$ as $\\epsilon\\to 0$, we obtain to\nleading order in ${1\/D}$ that $\\Delta U_0=0$ with $\\partial_{n} U_0=0$\non $\\partial\\Omega$. As such, we have that $U_0=U_0(t)$. At next\norder, we have that $U_1$ satisfies\n\\begin{equation}\\label{d:u1}\n\\Delta U_1=U_0+\\tau U_{0t}, \\quad \\vecb x \\in \\Omega\\backslash \\{\\vecb x_0\\}\\,;\n \\qquad \\partial_{\\vecb n} U_1 = 0, \\quad \\vecb x \\in \\partial\\Omega \\,.\n\\end{equation}\nThe formulation of this problem is complete once we determine a matching\ncondition for $U_1$ as $\\vecb x\\rightarrow \\vecb x_0$. \n\nTo obtain this matching condition, we must consider the inner region defined\nat ${\\mathcal O}(\\epsilon)$ distances outside the cell. In this inner\nregion we introduce the new variables \n$\\vecb y=\\epsilon^{-1}(\\vecb x-\\vecb x_0)$ and\n$\\hat{U}(\\vecb y,t) = U(\\vecb x_0+\\epsilon \\vecb y,t)$. From (\\ref{d:mainU}),\nwe obtain that\n\\begin{equation*}\n\\tau \\hat{U_t}=\\frac{D}{\\epsilon^2}\\Delta_{\\vecb y}\\hat{U}-\\hat{U}, \\quad\n \\rho=|\\vecb y|\\geq 1 \\,; \\qquad\nD\\frac{\\partial \\hat{U}}{\\partial \\rho} =d_1 \\hat{U}-d_2 u_1 \\,,\n \\quad \\mbox{on} \\quad \\rho=1 \\,.\n\\end{equation*}\nFor $D\\gg 1$, we then seek a radially symmetric solution to this inner\nproblem in the form\n\\begin{equation}\\label{d:uin}\n \\hat{U}(\\rho,t)=\\hat{U}_0(\\rho,t) +\\frac{1}{D}\\hat{U}_1(\\rho,t)+ \\cdots \\,.\n\\end{equation}\nTo leading order we obtain $\\Delta_{\\rho}\\hat{U}_0=0$ in $\\rho\\geq 1$,\nwith $\\hat{U}_{0\\rho}=0$ on $\\rho=1$, subject to the matching condition\nto the bulk that $\\hat{U}_0\\to U_0$ as $\\rho\\to \\infty$. We conclude\nthat $\\hat{U}_0=U_0$. At next order, the problem for $\\hat{U}_1$ is\n\\begin{equation}\\label{di:u1}\n\\Delta_{\\rho}\\hat{U}_1 = 0\\,, \\qquad \\rho\\ge 1\\,; \\qquad\n\\frac{\\partial \\hat{U}_1}{\\partial \\rho}=d_1 U_0-d_2 u_1\\,, \\qquad \\rho=1\\,.\n\\end{equation}\nAllowing for logarithmic growth at infinity, the solution to this problem is\n\\begin{equation}\\label{di:u1_sol}\n\\hat{U}_1=(d_1 U_0-d_2u_1)\\log\\rho+C \\,,\n\\end{equation}\nwhere $C$ is a constant to be found. Then, by writing (\\ref{di:u1_sol})\nin outer variables, and recalling (\\ref{d:uin}), we obtain that the \nfar-field behavior of the inner expansion is\n\\begin{equation}\\label{di:ff}\n \\hat{U}\\sim U_0 + \\frac{1}{D} \\left[ \\left(d_1U_0-d_2u_1\\right)\n \\log|\\vecb x-\\vecb x_0| + \\frac{1}{\\nu} \\left(d_1U_0-d_2u_1\\right) + C\\right]\n + \\cdots \\,.\n\\end{equation}\n\nFrom (\\ref{di:ff}) we observe that the term proportional to ${1\/D}$ is\nsmaller than the first term provided that $D\\gg O(\\nu^{-1})$. This\nspecifies the asymptotic range of $D$ for which our analysis will\nhold. From asymptotic matching of the bulk and inner solutions, the\nfar-field behavior of the inner solution (\\ref{di:ff}) provides the\nrequired singular behavior as $\\vecb x\\to \\vecb x_0$ for the outer bulk\nsolution. In this way, we find from (\\ref{di:ff}) and (\\ref{d:uout}) that\n$U_1$ satisfies (\\ref{d:u1}) subject to the singular behavior\n\\begin{equation}\\label{d:u1_sing}\n U_1 \\sim \\left(d_1U_0-d_2u_1\\right)\\log|\\vecb x - \\vecb x_0| \\,, \\quad\n \\mbox{as} \\quad \\vecb x\\to \\vecb x_0 \\,.\n\\end{equation}\nThen, (\\ref{d:u1}) together with (\\ref{d:u1_sing}) determines $U_1$ uniquely.\nFinally, in terms of this solution, we identify that the constant\n$C$ in (\\ref{di:ff}) and (\\ref{di:u1_sol}) is obtained from\n\\begin{equation}\\label{d1:ceval}\n\\lim_{\\vecb x\\to\\vecb x_0}\\left[U_1 -\n \\left(d_1U_0-d_2u_1\\right)\\log|\\vecb x - \\vecb x_0| \\right] =\n \\nu^{-1} \\left(d_1U_0-d_2u_1\\right) + C \\,.\n\\end{equation}\n\nWe now carry out the details of the analysis. We can write the problem\n(\\ref{d:u1}) and (\\ref{d:u1_sing}) for $U_1$ as\n\\begin{equation}\\label{d:u1_new}\n\\Delta U_1=U_0+\\tau U_{0t} + 2\\pi \\left(d_1U_0-d_2 u_1\\right) \n \\delta(\\vecb x-\\vecb x_0)\\,, \\qquad \\vecb x \\in \\Omega\\,; \n \\qquad \\partial_{n} U_1 = 0, \\quad \\vecb x \\in \\partial\\Omega \\,.\n\\end{equation}\nBy the divergence theorem, this problem has a solution only if\n$\\left(U_0+\\tau U_{0t}\\right)|\\Omega|=-2\\pi (d_1U_0-d_2 u_1)$. This leads to \nthe following ODE for the leading-order bulk solution $U_0(t)$:\n\\begin{equation}\\label{dlarge:one}\n U_0^{\\prime}= -\\frac{1}{\\tau} \\left( 1 + \\frac{2\\pi d_1}{|\\Omega|}\\right)\n U_0 + \\frac{2\\pi d_2}{\\tau |\\Omega|} u_1 \\,.\n\\end{equation}\nWithout loss of generality we impose $\\int_{\\Omega} U_1\\, d\\vecb x=0$\nso that $U_0$ is the spatial average of $U$. Then, the solution to\n(\\ref{d:u1_new}) is\n\\begin{equation}\\label{d:u1_sol}\n U_1 = -2\\pi \\left(d_1 U_0 - d_2 u_1\\right) G_0(\\vecb x;\\vecb x_0) \\,,\n\\end{equation}\nwhere $G_0(\\vecb x;\\vecb x_0)$ is the Neumann Green's function defined\nby (\\ref{g0:neum}). We then expand (\\ref{d:u1_sol}) as $\\vecb x\\to\n\\vecb x_0$, and use (\\ref{d1:ceval}) to identify $C$ in terms of\nthe regular part $R_0$ of the Neumann Green's function, defined in\n(\\ref{g0:neum}), as\n\\begin{equation}\\label{d:cres}\n C = - \\left(d_1 U_0 - d_2 u_1\\right) \\left(\\nu^{-1}+2\\pi R_0\\right)\n \\,.\n\\end{equation}\n\nIn summary, by using (\\ref{d:uin}), (\\ref{di:u1_sol}), and (\\ref{d:cres}),\nthe two-term inner expansion near the cell is given by\n\\begin{equation}\\label{ds:ui}\n \\hat{U}\\sim U_0 + \\frac{1}{D}\\left(d_1U_0-d_2u_1\\right)\\left(\n \\log\\rho - \\frac{1}{\\nu} - 2\\pi R_0 \\right) + \\cdots \\,.\n\\end{equation}\nFrom (\\ref{d:uout}) and (\\ref{d:u1_sol}), the two-term expansion\nfor the outer bulk solution, in terms of $U_0(t)$ satisfying the ODE\n(\\ref{dlarge:one}), is\n\\begin{equation}\n U \\sim U_0 - \\frac{2\\pi}{D} \\left(d_1 U_0 - d_2 u_1\\right) \n G_0(\\vecb x;\\vecb x_0) \\,. \\label{ds:uout}\n\\end{equation}\nThe final step in the analysis is to use (\\ref{d:mainuj}) to derive\nthe dynamics inside the cell. We readily calculate that\n\\begin{equation*}\n \\frac{1}{\\epsilon \\tau}\n \\int_{\\partial\\Omega_\\epsilon} \\left(d_1U - d_2u_1\\right)\\, ds\n \\sim \\frac{2\\pi}{ \\tau} \\left(d_1U_0- d_2u_1\\right) \\,,\n\\end{equation*}\nwhich determines the dynamics inside the cell from (\\ref{d:mainuj}).\n\nThis leads to our main result that, for $D\\gg {\\mathcal O}(\\nu^{-1})$,\nthe coupled PDE model (\\ref{d:mainbd}) reduces to the study of the\ncoupled (n+1) dimensional ODE system for the leading-order average\nbulk concentration $U_0(t)$ and cell dynamics $\\vecb u$ given by\n\\begin{equation}\\label{d:odes}\n U_0^{\\prime}= -\\frac{1}{\\tau} \\left( 1 + \\frac{2\\pi\n d_1}{|\\Omega|}\\right) U_0 + \\frac{2\\pi d_2}{\\tau |\\Omega|} u_1\n \\,, \\qquad \\vecb u^{\\prime} = \\vecb F(\\vecb u) + \\frac{2\\pi}{\n \\tau} \\left[d_1U_0-d_2u_1\\right] \\vecb e_1 \\,.\n\\end{equation}\n\nBefore studying (\\ref{d:odes}) for some specific reaction kinetics, we\nfirst examine conditions for the existence of steady-state solutions for\n(\\ref{d:odes}) and we derive the spectral problem characterizing the\nlinear stability of these steady-states.\n\nA steady-state solution $U_{0e}$ and $\\vecb u_e$ of (\\ref{d:odes}), if\nit exists, is a solution to the nonlinear algebraic system\n\\begin{equation} \\label{dlarge:ss}\n \\vecb F(\\vecb u_e) + \\frac{2\\pi}{ \\tau} \\left(d_1 U_{0e}-d_2\\vecb\n u_{1e}\\right) \\vecb e_1 = 0 \\,, \\qquad rU_{0e}= s u_{1e} \\,, \\qquad\n \\mbox{where} \\qquad r \\equiv 1 + \\frac{2\\pi d_1}{|\\Omega|} \\,, \\qquad \n s \\equiv \\frac{2\\pi d_2}{|\\Omega|}\\,.\n\\end{equation}\nTo examine the linearized stability of such a steady-state, we\nsubstitute $U_0=U_{0e}+e^{\\lambda t} \\eta$ and $\\vecb u =\\vecb u_e +\ne^{\\lambda t}\\vecb \\phi$ into (\\ref{d:odes}) and linearize. This\nyields that $\\eta$ and $\\vecb \\phi$ satisfy\n\\begin{equation*}\n \\lambda \\vecb \\phi = J \\vecb \\phi + \\frac{2\\pi}{ \\tau}\n \\left(d_1\\eta - d_2 \\phi_1\\right)\n \\vecb e_1 \\,, \\qquad \\tau \\lambda \\eta = - r\\eta + s\\phi_1 \\,,\n\\end{equation*}\nwhere $J$ is the Jacobian of $\\vecb F$ evaluated at the steady-state\n$\\vecb u=\\vecb u_e$. Upon solving the second equation for $\\eta$, and\nsubstituting the resulting expression into the first equation, we\nreadily derive the homogeneous linear system \n\\begin{equation}\\label{dlarge:eigv}\n \\left[(\\lambda I - J) - \\mu {\\mathcal E}_1\\right] \\vecb \\phi = 0\n \\,, \\qquad \\mbox{where} \\qquad \\mu \\equiv \\frac{2\\pi}{\n \\tau} \\left( \\frac{d_1 s}{\\tau\\lambda + r} - d_2\\right) \\,, \\qquad\n {\\mathcal E}_1 \\equiv \\vecb e_1\\vecb e_1^T \\,.\n\\end{equation}\n\nBy using the matrix determinant lemma we conclude that $\\lambda$ is an\neigenvalue of the linearization if and only if $\\lambda$ satisfies\n$\\vecb e_1^T \\left(\\lambda I - J\\right)^{-1}\\vecb e_1 = {1\/\\mu}$, where\n$\\mu$ is defined in (\\ref{dlarge:eigv}). From this expression, and by\nusing $d_1 s- d_2 r=-d_2$ as obtained from (\\ref{dlarge:ss}), we\nconclude that $\\lambda$ must be a root of ${\\mathcal Q}(\\lambda)=0$,\ndefined by\n\\begin{equation}\\label{dlarge:stab_crit}\n {\\mathcal Q}(\\lambda) \\equiv \n \\frac{{ \\tau} (r + \\tau\\lambda)}{2\\pi d_2 \\left(\n 1 + \\tau\\lambda\\right)} + \\frac{M_{11}}{\\det(\\lambda I - J)} \\,,\n\\end{equation}\nwhere $r$ is defined in (\\ref{dlarge:ss}). Here $M_{11}$\nis the cofactor of the element in the first row and first column of\n$\\lambda I-J$.\n\nNext, we show that (\\ref{dlarge:stab_crit}), which characterizes the\nstability of a steady-state solution of the ODE dynamics\n(\\ref{d:odes}), can also be derived by taking the limit $D_0\\gg 1$ in\nthe stability results obtained in (\\ref{dlarge:sync}) of \\S\n\\ref{sec:largeD} for the $D={\\mathcal O}(\\nu^{-1})$ regime where we\nset $D={D_0\/\\nu}$. Recall from the analysis in \\S \\ref{sec:largeD} for\n$D={D_0\/\\nu}$, that when $m=1$ only the synchronous mode can occur,\nand that the linearized eigenvalue satisfies\n(\\ref{dlarge:sync_new}). By formally letting $D_0\\to \\infty$ in\n(\\ref{dlarge:sync_new}) we readily recover (\\ref{dlarge:stab_crit}).\n\n\\subsection{Large D Theory: Analysis of Reduced Dynamics}\\label{dlarge:examples}\n\nWe now give some examples of our stability theory. We first consider\nthe case where there is exactly one dynamical species in the cell so\nthat $n=1$. From (\\ref{dlarge:ss}) with $n=1$ we obtain that the\nsteady-state $u_e$ is any solution of\n\\begin{equation}\\label{dlarge:1ode_ss}\n F(u_e) -\\frac{2\\pi d_2}{ \\tau} \\left[1+ \\frac{2\\pi\n d_1}{|\\Omega|}\\right]^{-1} u_e =0 \\,, \\qquad U_{0e}= \\frac{2\\pi\n d_2}{|\\Omega|} \\left[1+ \\frac{2\\pi d_1}{|\\Omega|}\\right]^{-1} u_e\n \\,.\n\\end{equation}\nIn the stability criterion (\\ref{dlarge:stab_crit}) we set $M_{11}=1$\nand $\\det(\\lambda I-J)= \\lambda-F_{u}^{e}$, where\n$F_{u}^{e}\\equiv {dF\/du}\\vert_{u=u_{e}}$, to conclude that the\nstability of this steady-state is determined by the roots of the\nquadratic \\begin{subequations} \\label{dlarge:stab_1cell}\n\\begin{equation}\\label{dlarge:quad}\n \\lambda^2 - \\lambda p_1 + p_2 =0 \\,,\n\\end{equation}\nwhere $p_1$ and $p_2$ are defined by\n\\begin{equation}\\label{dlarge:1cell_p12}\n p_1 = -\\frac{1}{\\tau} \\left(1+ \\frac{2\\pi d_1}{|\\Omega|}\\right) +\n F_u^{e} - \\frac{2\\pi d_2}{ \\tau} \\,, \\qquad p_2 =\n -\\frac{F_{u}^{e}}{\\tau} \\left(1+ \\frac{2\\pi\n d_1}{|\\Omega|}\\right) + \\frac{2\\pi d_2}{ \\tau^2} \\,.\n\\end{equation}\n\\end{subequations}$\\!$\nWe now establish the following result based on (\\ref{dlarge:stab_1cell}).\n\n\\begin{result}\\label{dlarge:theorem_1}\nLet $n=1$. Then, no steady-state solution of (\\ref{d:odes}) can undergo\na Hopf bifurcation. Moreover, if\n\\begin{equation}\\label{dlarge:stab_one}\n \\qquad F_{u}^{e} < F_{\\textrm{th}} \\equiv \\frac{2\\pi d_2}{ \\tau} \n \\left[ 1 + \\frac{2\\pi d_1}{|\\Omega|}\\right]^{-1} \\,,\n\\end{equation}\nthen $\\mbox{Re}(\\lambda)<0$, and so the steady-state is linearly stable.\nIf $F_{u}^{e} > F_{\\textrm{th}}$, the linearization has exactly one positive\neigenvalue.\n\\end{result}\n\n{\\ni \\bf \\underline{Proof:} } We first prove that no Hopf bifurcations are possible for the\nsteady-state. From (\\ref{dlarge:quad}) it is clear that there exists\na Hopf bifurcation if and only if $p_1=0$ and $p_2>0$ in\n(\\ref{dlarge:1cell_p12}). Upon setting $p_1=0$, we get\n$F_{u}^{e}=2\\pi d_2 { \\tau^{-1}} + \\tau^{-1}\\left(1+{2\\pi\n d_1\/|\\Omega|}\\right)$. Upon substituting this expression into\n(\\ref{dlarge:1cell_p12}) for $p_2$, we get that\n\\begin{equation*}\n p_2 = -\\frac{1}{\\tau} \\left[ \\frac{4\\pi^2 d_2^2}{{ \\tau}\n |\\Omega|} + \\frac{1}{\\tau}\n \\left( 1 + \\frac{2\\pi d_1}{|\\Omega|}\\right)\n \\left( 1 + \\frac{2\\pi d_2}{|\\Omega|}\\right)\\right]<0\\,.\n\\end{equation*}\nSince $p_2<0$ whenever $p_1=0$, we conclude that no Hopf bifurcations\nfor the steady-state are possible.\n\nThe second result follows by showing that $p_1<0$ and $p_2>0$\nwhen $F_{u}^{e} < F_{\\textrm{th}}$. From (\\ref{dlarge:1cell_p12}),\n$p_1<0$ and $p_2>0$ when\n\\begin{equation}\n \\frac{2\\pi d_2}{ \\tau} - F_u^{e} + \\frac{1}{\\tau} \\left(1+\n \\frac{2\\pi d_1}{|\\Omega|}\\right) >0 \\,, \\qquad \\frac{2\\pi\n d_2}{ \\tau} - F_u^{e} -\\frac{2\\pi d_1}{|\\Omega|}\n F_u^{e}>0 \\,.\n\\end{equation}\nThese two inequalities hold simultaneously only when the second relation\nis satisfied. This yields that when (\\ref{dlarge:stab_one}) holds, we have\n$\\mbox{Re}(\\lambda)<0$. Finally, when $F_{u}^{e} > F_{\\textrm{th}}$, we have\n$p_2<0$, and so there is a unique positive eigenvalue.\\hfill$\\blacksquare$\\vspace{6pt}\n\nThis result shows that with cell-bulk coupling the steady-state of the\nODE (\\ref{d:odes}) can be linearly stable even when the reaction\nkinetics is self-activating in the sense that $F_{u}^{e}>0$. Moreover,\nwe observe that as $\\tau$ increases, corresponding to when the\nintracellular kinetics has increasingly faster dynamics than the time\nscale for bulk decay, then the stability threshold $F_{\\textrm{th}}$\ndecreases. Therefore, for fast cell dynamics there is a smaller\nparameter range where self-activation of the intracelluar dynamics can\noccur while still maintaining stability of the steady-state to the\ncoupled system.\n\nNext, we let $n=2$ where $\\vecb F(\\vecb\nu)=(F(u_1,u_2),G(u_1,u_2))^T$. Then, the steady-state of the ODEs\n(\\ref{d:odes}) satisfies\n\\begin{equation}\\label{d:ss_2}\n F(u_{1e},u_{2e})-\\frac{2\\pi d_2}{r \\tau} u_{1e}=0 \\,, \\qquad\n G(u_{1e},u_{2e})=0 \\,, \\qquad U_{0e}=\\frac{s}{r} u_{1e} \\,,\n\\end{equation}\nwhere $r$ and $s$ are defined in (\\ref{dlarge:ss}). We then\nobserve that the roots of ${\\mathcal Q}(\\lambda)=0$ in\n(\\ref{dlarge:stab_crit}) are roots to a cubic polynomial in $\\lambda$.\nSince $M_{11}=\\lambda-G_{u_2}^{e}$, $\\det(\\lambda\nI-J)=\\lambda^2- \\mbox{tr}(J)\\lambda+\\det J$, where\n\\begin{equation}\n\\mbox{tr}(J)=F_{u_1}^e+G_{u_2}^e\\,, \\qquad \\det\nJ=F_{u_1}^eG_{u_2}^e-F_{u_2}^eG_{u_1}^e\\,,\n\\end{equation}\nand $F_{v}^e$, $G_{v}^e$ are partial derivatives of $F$ or $G$ with\nrespect to $v\\in(u_1,u_2)$ evaluated at the steady-state, we conclude\nthat the linear stability of the steady-state is characterized\nby the roots of the cubic\n\\begin{subequations} \\label{largeD:3ode}\n\\begin{equation}\\label{largeD:H3}\n\\lambda^3+p_1\\lambda^2+p_2\\lambda+p_3=0\\,,\n\\end{equation}\nwhere $p_1$, $p_2$ and $p_3$ are defined as\n\\begin{equation}\\label{largeD:p123}\n\\begin{split}\np_1&\\equiv \\frac{2\\pi d_2}{ \\tau}+\\frac{1}{\\tau}\\left(1+\\frac{2\\pi\n d_1}{|\\Omega|}\\right)-\\mbox{tr}(J)\\,, \\qquad\np_3\\equiv\\frac{1}{\\tau}\\left(\\left(1+\\frac{2\\pi d_1}{|\\Omega|}\\right)\\det J -\n \\frac{2\\pi d_2}{ \\tau} G_{u2}^e\\right)\\,, \\\\\n p_2 &\\equiv \\det J- \\frac{2\\pi\nd_2}{ \\tau} G_{u_2}^e+\\frac{1}{\\tau}\\left(\n \\frac{2\\pi d_2}{ \\tau}\n - \\left(1+\\frac{2\\pi\n d_1}{|\\Omega|}\\right) \\mbox{tr}(J) \\right)\\,.\n\\end{split}\n\\end{equation}\n\\end{subequations}$\\!$\nBy taking $m=1$ and letting $D_0\\rightarrow \\infty$ in (\\ref{mcell2ode:p123}),\nit is readily verified that the expressions above for $p_i$, for $i=1,2,3$,\nagree exactly with those in (\\ref{mcell2ode:p123}). Then, by satisfying\nthe Routh-Hurwitz conditions (\\ref{mcell2ode:syn}), we can\nplot the Hopf bifurcation boundaries in different parameter planes.\n\n\\subsubsection{Example: One Cell with Sel'kov Dynamics}\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{d1d2_largeDcomp.eps}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{ode_sim.eps}\n\\caption{Left panel: Comparison of the Hopf bifurcation boundaries for\n the synchronous mode for the Sel'kov model (\\ref{2dsel:fg}) in the\n $d_1$ versus $d_2$ plane with $D_0=1,\\ 5,\\ 50$ (solid), as\n obtained from (\\ref{d0:cubic}), and the large $D$ approximation\n (dashed), as obtained from (\\ref{largeD:3ode}). Between the outer\n two black curves, the synchronous mode is unstable for $D_0=1$.\n In the region bounded by the solid or dashed curve this\n mode is unstable. As $D_0$ increases, the Hopf\n boundaries from (\\ref{d0:cubic}) gradually approach the one given in\n from the large $D$ approximation. Right panel:\n Numerical results for the ODE system (\\ref{d:odes}) showing\n sustained oscillations. In the left and right panels we fixed\n $\\mu=2$, $\\alpha=0.9$, $\\epsilon_0=0.15$, $\\tau=1$, and\n $|\\Omega|=10$. In the right panel we took $d_1=0.8$ and\n $d_2=0.05$ for which the steady-state solution of the ODEs\n (\\ref{d:odes}) is unstable.}\n\\label{dlarge:hopfbd_sel}\n\\end{center}\n\\end{figure}\n\nWhen the intracellular kinetics is described by the Sel'kov model,\nwhere $F$ and $G$ are given in (\\ref{2dsel:fg}), we obtain from\n(\\ref{d:ss_2}) that the steady-state solution of the ODE dynamics is\n\\begin{equation}\nu_1^e=\\frac{r\\mu}{\\left[r+{2\\pi d_2\/ \\tau} \\right]}\\,,\\qquad \nu_2^e=\\frac{\\mu}{\\alpha+(u_1^e)^2}\\,,\\qquad U_{0e}=\n\\frac{s\\mu}{r+{2\\pi d_2\/\\tau}}\\,,\n\\end{equation}\nwhere $r$ and $s$ are defined in (\\ref{dlarge:ss}). The partial\nderivatives of $F$ and $G$ can be calculated as in\n(\\ref{2dsel:fgp}). \n\nIn the left panel of Fig.~\\ref{dlarge:hopfbd_sel} we plot the Hopf\nbifurcation boundary in the $d_1$ versus $d_2$ plane associated with\nlinearizing the ODEs (\\ref{d:odes}) about this steady-state solution.\nThis curve was obtained by setting $p_1 p_2=p_3$ with $p_1>0$ and\n$p_3>0$ in (\\ref{largeD:3ode}). In this figure we also plot the Hopf\nbifurcation boundary for different values of $D_0$, with\n${D=D_0\/\\nu}$, as obtained from our stability formulation\n(\\ref{d0:cubic}) of \\S \\ref{sec:examples} for the $D={\\mathcal\n O}(\\nu^{-1})$ regime. We observe from this figure that the stability\nboundary with $D_0=50$ closely approximates that obtained from\n(\\ref{largeD:3ode}), which corresponds to the $D_0\\to \\infty$ limit of\n(\\ref{d0:cubic}).\n\nWe emphasize that since in the distinguished limit $D\\gg{\\mathcal\n O}(\\nu^{-1})$ we can approximate the original coupled PDE system\n(\\ref{d:mainbd}) by the system (\\ref{d:odes}) of ODEs, the study\nof large-scale time dynamics far from the steady-state becomes\ntractable. In the right panel of Fig.~\\ref{dlarge:hopfbd_sel}, we plot\nthe numerical solution to (\\ref{d:odes}) with Sel'kov dynamics\n(\\ref{2dsel:fg}), where the initial condition is $u_1(0)=0.01$,\n$u_2(0)=0.2$, and $U_0(0)=0.5$. We observe that by choosing $d_1$ and\n$d_2$ inside the region bounded by the dashed curve in the left panel\nof Fig.~\\ref{dlarge:hopfbd_sel}, where the steady-state is unstable,\nthe full ODE system (\\ref{d:odes}) exhibits a stable periodic orbit,\nindicating a limit cycle behavior.\n\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{u12U_t.eps}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{u12_nocoupling.eps}\n\\caption{Left panel: $u_1$, $u_2$ and $U_0$ versus $t$ showing\n sustained oscillatory dynamics. Parameters are $\\mu=2$,\n $\\alpha=0.9$, $\\epsilon_0=0.15$, $\\tau=1$, $|\\Omega|=10$, $d_1=0.8$\n and $d_2=0.05$. Right panel: $u_2$ versus $u_1$ when the local\n kinetics is decoupled from the bulk, showing decaying oscillations\n towards the stable steady-state solution $u_1=\\mu$ and\n $u_2=\\mu\/(\\alpha+u_1^2)$. Initial values are $u_1(0)=0.01$ and\n $u_2(0)=0.2$. The parameter values of $\\mu$, $\\epsilon_0$ and $\\alpha$\n are the same as in the left panel.}\n\\label{dlarge:u12U}\n\\end{center}\n\\end{figure}\n\nIn addition, in the left panel of Fig.~\\ref{dlarge:u12U} we plot the\ntime evolution of $u_1$, $u_2$ and $U_0$, showing clearly the\nsustained periodic oscillations. For comparison, fixing the same\nparameter set for the Sel'kov kinetics (\\ref{2dsel:fg}), in the right\npanel of Fig.~\\ref{dlarge:u12U} we plot the phase plane of $u_2$\nversus $u_1$ when there is no coupling between the local kinetics and\nthe bulk. We now observe that without this cell-bulk coupling the\nSel'kov model (\\ref{2dsel:fg}) exhibits transient decaying oscillatory\ndynamics, with a spiral behavior in the phase-plane towards the\nlinearly stable steady-state.\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{bif_sel_largeD.eps}\n\\caption{Global bifurcation diagram of $u_1$ with respect to $d_2$ at\n a fixed $d_1=0.8$ as computed using XPPAUT \\cite{xpp} from the ODE\n system (\\ref{d:odes}) for the Sel'kov kinetics (\\ref{2dsel:fg}). The\n thick\/thin solid line represents stable\/unstable steady-state\n solutions of $u_1$, while the solid dots indicate a stable periodic\n solution branch. The parameters used are $\\mu=2$, $\\alpha=0.9$,\n $\\epsilon_0=0.15$, $\\tau=1$, and $|\\Omega|=10$. }\\label{sel:xpp}\n\\end{center}\n\\end{figure}\n\nFinally, we use the numerical bifurcation software XPPAUT \\cite{xpp}\nto confirm the existence of a stable large amplitude periodic solution\nto (\\ref{d:odes}) with Sel'kov kinetics when $d_1$ and $d_2$ are in the\nunstable region of the left panel of Fig.~\\ref{dlarge:hopfbd_sel}. In\nFig.~\\ref{sel:xpp} we plot a global bifurcation diagram of $u_{1}$\nversus $d_2$ for $d_1=0.8$, corresponding to taking a horizontal\nslice at fixed $d_2=0.8$ through the stability boundaries in the $d_2$\nversus $d_1$ plane shown in Fig.~\\ref{dlarge:hopfbd_sel}. The two\ncomputed Hopf bifurcation points at $d_2\\approx 0.0398$ and\n$d_1=0.1098$ agree with the theoretically predicted values in\nFig.~\\ref{dlarge:hopfbd_sel}. \n\n\\subsubsection{Example: One Cell with Fitzhugh-Nagumo Dynamics}\n\nFinally, we apply our large $D$ theory to the case where the\nintracellular dynamics is governed by the FN kinetics\n(\\ref{2dsf:fg}). From (\\ref{d:ss_2}) we obtain that the steady-state\nsolution of the ODEs (\\ref{d:odes}) with the kinetics (\\ref{2dsf:fg})\nis\n\\begin{equation}\\label{dlarge:fn_eq}\nu_1^e=\\Lambda u_{2}^e \\,, \\qquad U_{0e}=\\frac{s u_{1}^e}{r} \\,, \n \\qquad \\mbox{where} \\qquad\n \\Lambda \\equiv \\frac{\\epsilon_0 z r}{\\left[\\epsilon_0 r + {2\\pi d_2\/{\n \\tau}}\\right]} \\,.\n\\end{equation}\nHere $r$ and $s$ are defined in (\\ref{dlarge:ss}), and $u_{2e}>0$\nis the unique root of the cubic (\\ref{2dfitz:cubic}) where $\\Lambda$\nin (\\ref{2dfitz:cubic}) is now defined in (\\ref{dlarge:fn_eq}). The\npartial derivatives of $F$ and $G$ can be calculated as in\n(\\ref{2dsf:jac}).\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{d1d2_slowfast_largeDcomp.eps}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{ode_sim_slowfast.eps}\n\\caption{Left panel: Comparison of the Hopf bifurcation boundaries for\n the synchronous mode with FN kinetics (\\ref{2dsf:fg})\n in the $d_1$ versus $d_2$ parameter plane with $D_0=1,\\ 5,\\ 50$\n (solid), as obtained from (\\ref{d0:cubic}), and the large $D$\n approximation (dashed), as obtained from (\\ref{largeD:3ode}). In the\n wedge-shaped regions bounded by the solid curves the synchronous\n mode is unstable for the finite values of $D_0$. As $D_0$ increases,\n the Hopf boundaries obtained from (\\ref{d0:cubic}) becomes rather close\n to the dashed one obtained from (\\ref{largeD:3ode}) from the large\n $D$ approximation. Right panel: Numerical simulation for the ODE\n system (\\ref{d:odes}), showing sustained oscillations, with initial\n conditions $u_1(0)=6.0$, $u_2(0)=2.3$, and $U_0(0)=0.05$. In the left\n and right panels we fixed $z=3.5$, $q=5$, $\\epsilon_0=0.5$,\n $\\tau=1$, and $|\\Omega|=10$, and in the right panel we took\n $d_1=15.6$ and $d_2=0.2$ corresponding to a point where the\n steady-state solution of the ODEs (\\ref{d:odes}) is unstable.}\n\\label{dlarge:hopfbd_fn}\n\\end{center}\n\\end{figure}\n\nIn the left panel of Fig.~\\ref{dlarge:hopfbd_fn} the dashed curve is\nthe Hopf bifurcation boundary in the $d_1$ versus $d_2$ plane\nassociated with linearizing the ODEs (\\ref{d:odes}) about this\nsteady-state solution. In this figure the Hopf bifurcation boundaries\nfor different values of $D_0$, with ${D=D_0\/\\nu}$, are also\nshown. These latter curves were obtained from our stability\nformulation (\\ref{d0:cubic}) of \\S \\ref{sec:examples}. Similar to what\nwe found for the Sel'kov model, the stability boundary for $D_0=50$ is\nvery close to that for the infinite $D$ result obtained from\n(\\ref{largeD:3ode}). In the right panel of Fig.~\\ref{dlarge:hopfbd_fn}\nwe plot the numerical solution to (\\ref{d:odes}) with FN dynamics\n(\\ref{2dsf:fg}) for the parameter set $d_1=15.6$ and $d_2=0.2$, which\nis inside the unstable region bounded by the dashed curves in the left\npanel of Fig.~\\ref{dlarge:hopfbd_fn}. With the initial condition\n$u_1(0)=6.0$, $u_2(0)=2.3$, and $U_0(0)=0.05$, the numerical\ncomputations of the full ODE system (\\ref{d:odes}) again reveal a\nsustained and stable periodic solution.\n\nFinally, we use XPPAUT \\cite{xpp} on (\\ref{d:odes}) to compute a\nglobal bifurcation of $u_1$ versus $d_1$ for fixed $d_2=0.2$ for FN\nkinetics. This plot corresponds to taking a vertical slice at fixed\n$d_2=0.2$ through the stability boundaries in the $d_1$ versus $d_2$\nplane shown in Fig.~\\ref{dlarge:hopfbd_fn}. The two computed Hopf\nbifurcation points at $d_1\\approx 15.389$ and $d_1=42.842$ agree with\nthe predicted values in Fig.~\\ref{dlarge:hopfbd_fn}. These results\nconfirm the existence of a stable periodic solution branch induced by\nthe cell-bulk coupling.\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{bif_slowfast_largeD.eps}\n\\caption{Global bifurcation diagram of $u_{1}$ versus $d_1$ when\n $d_2=0.2$, as computed using XPPAUT \\cite{xpp} from the ODE system\n (\\ref{d:odes}) for the FN kinetics (\\ref{2dsf:fg}). The thick\/thin\n solid line represents stable\/unstable steady-state solutions,\n while the solid dots indicate a stable periodic solution.\n The other parameters are $z=3.5$, $q=5$, $\\epsilon_0=0.5$,\n $\\tau=1$, and $|\\Omega|=10$. }\n\\label{fn:xpp}\n\\end{center}\n\\end{figure}\n\n\\setcounter{equation}{0}\n\\setcounter{section}{5}\n\\section{The Effect of the Spatial Configuration of the Small Cells: \nThe $D={\\mathcal O}(1)$ Regime}\\label{sec:finite_d}\n\nIn this section we construct steady-state solutions and study their\nlinear stability properties in the $D={\\mathcal O}(1)$ regime, where\nboth the number of cells and their spatial distribution in the domain\nare important factors. For simplicity, we consider a special spatial\nconfiguration of the cells inside the unit disk $\\Omega$ for which the\nGreen's matrix ${\\mathcal G}$ has a cyclic structure. More\nspecifically, on a ring of radius $r_0$, with $00$ leads to an instability of the steady-state\nsolution on an ${\\mathcal O}(1)$ time-scale. If all such roots satisfy\n$\\mbox{Re}(\\lambda)<0$, then the steady-state is linearly stable on an\n${\\mathcal O}(1)$ time-scale.\n\nTo study the stability properties of the steady-state using\n(\\ref{cyclic:omega_j}), and identify any possible Hopf bifurcation\nvalues, we must first calculate the spectrum (\\ref{cyclic:Glam_mat})\nof the cyclic and symmetric matrix ${\\mathcal G}_{\\lambda}$, whose\nentries are determined by the $\\lambda$-dependent reduced-wave Green's\nfunction $G_{\\lambda}(\\vecb x;\\vecb \\xi)$, with regular part\n$R_{\\lambda}(\\vecb \\xi)$, as defined by (\\ref{bdgneig}). Since\n${\\mathcal G}_{\\lambda}$ is not a Hermitian matrix when $\\lambda$ is\ncomplex, its eigenvalues $\\omega_{\\lambda,j}$ are in general\ncomplex-valued when $\\lambda$ is complex. Then, by replacing\n$\\theta_0$ in (\\ref{cyclic:G}) and (\\ref{cyclic:R}) with\n$\\theta_\\lambda\\equiv \\sqrt{(1+\\tau\\lambda)\/D}$, we readily obtain\nthat\n\\begin{equation} \\label{cyclic:Glam}\nG_\\lambda(\\vecb x;\\vecb \\xi) = \\frac{\\,1}{\\,2 \\pi}\nK_0\\left(\\theta_\\lambda|\\vecb x-\\vecb \\xi|\\right) -\\frac{\\,1}{\\,2 \\pi}\n\\sum^{\\infty}_{n=0} \\sigma_n \\cos\\left(n(\\psi-\\psi_0)\\right)\n\\frac{K^{\\prime}_n(\\theta_\\lambda)} {I^{\\prime}_n(\\theta_\\lambda)}\nI_n\\left(\\theta_\\lambda r\\right) I_n\\left(\\theta_\\lambda r_0\\right)\n\\,; \\quad \\sigma_0=1 \\,, \\quad \\sigma_n=2\\,, \\,\\,\\, n\\geq 1\\,,\n\\end{equation}\nwith regular part\n\\begin{equation}\\label{cyclic:Rlam}\n R_{\\lambda} = \\frac{1}{2\\pi} \\left[\\log{2} - \\gamma_e +\n \\frac{1}{2}\\log{D} -\\frac{1}{2}\\log(1+\\tau\\lambda) \\right]\n -\\frac{\\,1}{\\,2 \\pi} \\sum^{\\infty}_{n=0} \\sigma_n\n \\frac{K^{\\prime}_n(\\theta_\\lambda)} {I^{\\prime}_n(\\theta_\\lambda)} \\left[\n I_n\\left(\\theta_\\lambda r_0\\right) \\right]^2 \\,,\n\\end{equation}\nwhere we have specified the principal branch for $\\theta_\\lambda$.\nThe Green's matrix ${\\mathcal G}_\\lambda$ is obtained by a cyclic\npermutation of its first row $\\vecb a_\\lambda \\equiv\n(a_{\\lambda,1},\\ldots,a_{\\lambda,m})^T$, which is defined term-wise by\n\\begin{equation}\\label{cyclic:G_lam_row}\n a_{\\lambda,1} \\equiv R_{\\lambda} \\,; \\qquad a_{\\lambda,j} =\n G_{\\lambda,j1} \\equiv G_{\\lambda} (\\vecb x_j;\\vecb x_1) \\,, \\quad\n j=2,\\ldots, m \\,.\n\\end{equation}\nAgain we can numerically evaluate $G_{\\lambda,j1}$ for $j=2,\\ldots,m$\nand $R_\\lambda$ by using (\\ref{cyclic:Glam}) and (\\ref{cyclic:Rlam}),\nrespectively. \n\nNext, we must determine the full spectrum\n(\\ref{cyclic:Glam_mat}) of the cyclic and symmetric matrix ${\\mathcal\n G}_\\lambda$. For the $m\\times m$ cyclic matrix ${\\mathcal\n G}_{\\lambda}$, generated by permutations of the row vector $\\vecb\na_{\\lambda}$, it is well-known that its eigenvectors $\\vecb v_j$ and\neigenvalues $\\omega_{\\lambda,j}$ are\n\\begin{equation}\\label{cyclic:mat_eig}\n\\omega_{\\lambda, j}= \\sum_{n=0}^{m-1} a_{\\lambda, n+1} e^{2 \\pi i (j-1)n\/m} \\,,\n \\qquad \\vecb v_j = (1, e^{2 \\pi i (j-1)\/m} \\,, \\ldots, e^{2\n \\pi i (j-1) (m-1)\/m})^T \\,, \\qquad j=1,\\ldots, m\\,.\n\\end{equation} \nSince ${\\mathcal G}$ is also necessarily a symmetric matrix it\nfollows that\n\\begin{equation}\\label{cyclic:arow}\n a_{\\lambda,j} = a_{\\lambda,m+2-j} \\,, \\qquad j=2,\\ldots, \n\\lceil {m\/2} \\rceil \\,,\n\\end{equation}\nwhere the ceiling function $\\lceil x \\rceil$ is defined as the\nsmallest integer not less than $x$. This relation can be used to\nsimplify the expression (\\ref{cyclic:mat_eig}) for\n$\\omega_{\\lambda,j}$, into the form as written below in\n(\\ref{cyclic:spectrum_all}). Moreover, as a result of\n(\\ref{cyclic:arow}), it readily follows that\n\\begin{equation}\\label{cyclic:omega}\n \\omega_{\\lambda,j}=\\omega_{\\lambda,m+2-j} \\,, \\qquad \n \\mbox{for} \\quad j=2,\\ldots, \\lceil {m\/2} \\rceil \\,,\n\\end{equation}\nso that there are $\\lceil {m\/2} \\rceil -1$ eigenvalues of multiplicity\ntwo. For these multiple eigenvalues the two independent real-valued\neigenfunctions are readily seen to be $\\mbox{Re}(\\vecb v_j)={(\\vecb\n v_j +\\vecb v_{m+2-j})\/2}$ and $\\mbox{Im}(\\vecb v_j)={(\\vecb v_j\n -\\vecb v_{m+2-j})\/(2i)}$. In addition to $\\omega_1$, we also observe\nthat there is an additional eigenvalue of multiplicity one when $m$ is\neven.\n\nIn this way, our result for the matrix spectrum of ${\\mathcal G}_{\\lambda}$\nis as follows: The synchronous eigenpair of ${\\mathcal G}_{\\lambda}$ is \n\\begin{subequations} \\label{cyclic:spectrum_all}\n\\begin{equation}\\label{cyclic:spectrum_1}\n\\omega_{\\lambda, 1} = \\sum_{n=1}^m a_{\\lambda,n}\\,, \\qquad v_1 =\n(1, \\ldots, 1)^T\\,,\n\\end{equation}\nwhile the other eigenvalues, corresponding to the asynchronous modes, are\n\\begin{equation}\\label{cyclic:spectrum_j}\n\\omega_{\\lambda, j} = \\sum_{n=0}^{m-1} \\cos\\left( \\frac{2 \\pi\n(j-1)\\,n}{m}\\right) a_{\\lambda, n+1}\\,, \\qquad j=2,\\ldots,m \\,,\n\\end{equation}\nwhere $\\omega_{\\lambda,j}=\\omega_{\\lambda,m+2-j}$ for $j=2,\\ldots,\n\\lceil {m\/2} \\rceil$. When $m$ is even, we notice that there is an\neigenvalue of multiplicity one given by\n$\\omega_{\\lambda,\\frac{m}{2}+1}=\\sum_{n=0}^{m-1} (-1)^{n}\na_{n+1}$. The corresponding eigenvectors for\n$j=2,\\ldots, \\lceil {m\/2} \\rceil$ can be written as\n\\begin{equation}\\label{cyclic:eigvec}\n\\begin{aligned}\n \\vecb v_j &= \\left(1, \\cos \\left(\\frac{2 \\pi (j-1)}{m}\\right), \\ldots,\n \\cos \\left( \\frac{2 \\pi (j-1)(m-1)}{m} \\right)\\,\\right)^T \\,,\\\\\n \\vecb v_{m+2-j} &= \\left(0, \\sin \\left( \\frac{2 \\pi (j-1)}{m}\\right),\n \\ldots, \\sin \\left( \\frac{2 \\pi (j-1)(m-1)}{m}\\right)\\,\\right)^T \\,.\n\\end{aligned}\n\\end{equation}\n\\end{subequations}$\\!$\nFinally, when $m$ is even, there is an additional eigenvector given by\n$\\vecb v_{\\frac{m}{2}+1}=(1,-1,\\ldots,-1)^T$.\n\nWith the eigenvalues $\\omega_{\\lambda,j}$, for $j=1,\\ldots,m$,\ndetermined in this way, any Hopf bifurcation boundary in parameter\nspace is obtained by substituting $\\lambda=i\\lambda_I$ with\n$\\lambda_I>0$ into (\\ref{cyclic:omega_j}), and requiring that the real\nand imaginary parts of the resulting expression vanish. This yields,\nfor each $j=1,\\ldots,m$, that\n\\begin{equation}\\label{cyclic:real_imag}\n \\mbox{Re}\\left(\\omega_{\\lambda,j}\\right) + \\frac{1}{2\\pi\\nu} \\left(\n 1+ \\frac{D\\nu}{d_1} \\right) + \\frac{d_2 D} {d_1\\tau} \\mbox{Re}\\left(\n \\frac{M_{11}}{\\det(\\lambda I-J) }\\right)=0 \\,, \\qquad\n \\mbox{Im}\\left(\\omega_{\\lambda,j}\\right) + \\frac{d_2 D} {d_1\\tau}\n \\mbox{Im}\\left( \\frac{M_{11}}{\\det(\\lambda I-J) }\\right)=0 \\,.\n\\end{equation}\nFinally, we can use the winding number criterion of complex analysis\non (\\ref{cyclic:omega_j}) to count the number of eigenvalues of the\nlinearization when the parameters are off any Hopf bifurcation\nboundary. This criterion is formulated below in \\S \\ref{cyc:sel}.\n\nWe remark that in the limit $D\\gg 1$, we can use $K_{0}(z)\\sim\n-\\log{z}$ together with $I_{0}(z)\\sim 1+{z^2\/4}$ as $z\\to 0$, to\nestimate from the $n=0$ term in (\\ref{cyclic:Glam}) and\n(\\ref{cyclic:Rlam}) that\n$-(2\\pi)^{-1}{K_0^{\\prime}(\\theta_\\lambda)\/I_0^{\\prime}(\\theta_\\lambda)}\\sim\n{D\/\\left[\\pi(1+\\tau\\lambda)\\right]}$ as $D\\to \\infty$. Therefore, for\n$D\\to \\infty$, the Green's matrix ${\\cal G}_{\\lambda}$ satisfies\n${\\mathcal G}_\\lambda \\to {Dm {\\mathcal\n E}\/{\\left[\\pi(1+\\tau\\lambda)\\right]}}$, where ${\\mathcal\n E}={{\\vecb e}{\\vecb e}^T\/m}$ and $\\vecb e\\equiv\n(1,\\ldots,1)^T$. This yields for $D\\gg 1$ that\n$\\omega_1={Dm\/{\\left[\\pi (1+\\tau\\lambda)\\right]}}$ and\n$\\omega_j={\\mathcal O}(1)$ for $j=2,\\dots,n$. By substituting these\nexpressions into (\\ref{cyclic:real_imag}), we can readily recover the\nspectral problems (\\ref{dlarge:sync_new}) and (\\ref{dlarge:async}),\nconsidered in \\S \\ref{sec:largeD}, associated with the regime\n$D={\\mathcal O}(\\nu^{-1})$. Therefore, (\\ref{cyclic:real_imag})\nprovides a smooth transition to the leading-order spectral problems\nconsidered in \\S \\ref{sec:largeD} for $D={\\mathcal O}(\\nu^{-1})$.\n\n\\subsection{Example: The Sel'kov Model}\\label{cyc:sel}\n\nWe now use (\\ref{cyclic:real_imag}) to compute phase diagrams in the\n$\\tau$ versus $D$ parameter space associated with $m$ equally-spaced\ncells of radius $\\epsilon$ on a ring of radius $r_0$, with $00 \\,, \\qquad\n \\mbox{tr}(J)= \\frac{1}{\\alpha+u_{1e}^2}\\left[ 2 u_{1e}\\mu -\n (\\alpha+u_{1e}^2) - \\epsilon_0 (\\alpha+u_{1e}^2)^2\\right] \\,.\n\\end{equation}\n\\end{subequations}$\\!$ In this subsection we fix the Sel'kov parameters $\\mu$, $\\alpha$, and\n$\\epsilon_0$, the permeabilities $d_1$ and $d_2$, and the cell radius\n$\\epsilon$ as\n\\begin{equation}\n \\mu=2 \\,, \\qquad \\alpha=0.9\\,, \\qquad \\epsilon_0=0.15, \\qquad d_1=0.8\\,, \n \\qquad d_2=0.2 \\,, \\qquad \\epsilon=0.05 \\,. \\label{cyc:par}\n\\end{equation}\nWith these values for $\\mu$, $\\alpha$, and $\\epsilon_0$, the\nintracellular dynamics has a stable steady-state when uncoupled from\nthe bulk.\n\nThen, to determine the Hopf bifurcation boundary for the coupled\ncell-bulk model we set $M_{11}=\\lambda-G_{u_{2}}^e$ in\n(\\ref{cyclic:real_imag}), and use $G_{u_{2}}^e=-\\mbox{det}(J)$ as\nobtained from (\\ref{2dsel:jac}). By letting $\\lambda=i\\lambda_I$\nin the resulting expression, we conclude that any Hopf bifurcation\nboundary, for each mode $j=1,\\ldots,m$, must satisfy\n\\begin{equation}\\label{cyclic:hb_ex}\n \\begin{split}\n \\mbox{Re}\\left(\\omega_{\\lambda,j}\\right) + \\frac{1}{2\\pi\\nu} \\left(\n 1+ \\frac{D\\nu}{d_1} \\right) - \\left(\\frac{d_2 D} {d_1\\tau}\\right) \\frac{\n \\left[ \\lambda_I^2 \\mbox{tr}(J) + \\mbox{det}(J) \\left(\\lambda_I^2-\n \\mbox{det}(J)\\right)\\right]}{ \\left[ \\left(\\mbox{det}(J)-\\lambda_I^2\\right)^2\n + \\left(\\lambda_I \\mbox{tr}(J)\\right)^2 \\right]} =0 \\,,\\\\\n \\mbox{Im}\\left(\\omega_{\\lambda,j}\\right) + \\left(\\frac{d_2 D} {d_1\\tau}\\right)\n\\frac{ \\left[ \\lambda_I \\left(\\mbox{det}(J)-\\lambda_I^2\\right) + \n \\mbox{det}(J) \\mbox{tr}(J) \\lambda_I \\right]}\n { \\left[ \\left(\\mbox{det}(J)-\\lambda_I^2\\right)^2\n + \\left(\\lambda_I \\mbox{tr}(J)\\right)^2 \\right]} =0 \\,.\n \\end{split}\n\\end{equation}\nFor a specified value of $D$, we view (\\ref{cyclic:hb_ex}) as a\ncoupled system for the Hopf bifurcation value of $\\tau$ and the\ncorresponding eigenvalue $\\lambda_I$, which we solve by Newton's\nmethod. \n\nFor parameter values off of any Hopf bifurcation boundary, we\ncan use the winding number criterion on ${\\mathcal F}_j(\\lambda)$ in\n(\\ref{cyclic:omega_j}) to count the number of unstable eigenvalues\n$N_j$ of the linearization for the $j$-th mode. By using the\nargument principle, we obtain that the number $N_j$ of roots of\n${\\mathcal F}_j(\\lambda)=0$ in $\\mbox{Re}(\\lambda)>0$ is\n\\begin{equation}\\label{wind:last_1}\nN_j=\\frac{1}{2\\pi}[\\mbox{arg} {\\mathcal F}_j]_{\\Gamma}+P\\,, \n\\end{equation}\nwhere $P$ is the number of poles of ${\\mathcal F}_{j}(\\lambda)$ in\n$\\mbox{Re}(\\lambda)>0$, and the square bracket denotes the change in\nthe argument of ${\\mathcal F}_{j}$ over the contour $\\Gamma$. Here the\nclosed contour $\\Gamma$ is the limit as ${\\mathcal R}\\to \\infty$ of\nthe union of the imaginary axis, which can be decomposed as\n$\\Gamma_{I+}=i\\lambda_I$ and $\\Gamma_{I-}=-i\\lambda_I$, for\n$0<\\lambda_I<{\\mathcal R}$, and the semi-circle $\\Gamma_{\\mathcal R}$\ndefined by $|\\lambda|={\\mathcal R}$ with\n$|\\mbox{arg}(\\lambda)|\\leq\\pi\/2$. Since $\\omega_{\\lambda j}$ is\nanalytic in $\\mbox{Re}(\\lambda)>0$, it follows that $P$ is determined\nby the number of roots of $\\det(\\lambda I-J)=0$ in\n$\\mbox{Re}(\\lambda)>0$. Since $\\mbox{det}(J)>0$, as shown in\n(\\ref{cyc:det_tr}), we have that $P=2$ if $\\mbox{tr}(J)>0$ and $P=0$\nif $\\mbox{tr}(J)<0$. Next, we let ${\\mathcal R}\\to \\infty$ on\n$\\Gamma_{\\mathcal R}$ and calculate $[\\mbox{arg}\n \\mathcal{F}_{j}]_{\\Gamma_{\\mathcal R}}$. It is readily seen that the\nGreen's matrix ${\\mathcal G}_{\\lambda}$ tends to a multiple of a\ndiagonal matrix on $\\Gamma_{\\mathcal R}$ as ${\\mathcal R}\\gg 1$, of\nthe form ${\\mathcal G}_{\\lambda}\\to R_{\\lambda,\\infty}I$, where\n$R_{\\lambda,\\infty}$ is the regular part of the free-space Green's\nfunction $G_{f}(\\vecb x;\\vecb x_0) =\n(2\\pi)^{-1} K_{0}\\left(\\theta_\\lambda|\\vecb x-\\vecb x_0| \\right)$ at\n$\\vecb x=\\vecb x_0$, given explicitly by the first term in the\nexpression (\\ref{cyclic:Rlam}) for $R_{\\lambda}$. Since\n$\\omega_{\\lambda,j}\\to R_{\\lambda,\\infty}$ for $j=1,\\ldots,m$, we estimate\n on $\\Gamma_{\\mathcal R}$ as ${\\mathcal R}\\gg 1$ that\n\\begin{equation*}\n {\\mathcal F}_{j}(\\lambda) \\sim -\\frac{1}{2\\pi}\n \\log\\sqrt{1+\\tau\\lambda} + c_0 + {\\mathcal\n O}\\left({1\/\\lambda}\\right) \\,,\n\\end{equation*}\nfor some constant $c_0$. It follows that $\\mathcal{F}_j(\\lambda)\\sim\nO(\\ln {\\mathcal R})- {i\/8}$ as ${\\mathcal R}\\to \\infty$, so that\n$\\lim_{{\\mathcal R}\\to\n \\infty}[\\mbox{arg}\\mathcal{F}_{j}]_{\\Gamma_{\\mathcal R}}=0$.\nFinally, since $[\\mbox{arg}\\mathcal{F}_{j}]_{\\Gamma_{I+}}= [\\mbox{arg}\n \\mathcal{F}_{j}]_{\\Gamma_{I-}}$, as a consequence of ${\\mathcal F}_j$\nbeing real-valued when $\\lambda$ is real, we conclude from\n(\\ref{wind:last_1}) that\n\\begin{equation}\\label{wind:last_2}\nN_j=\\frac{1}{2\\pi}[\\mbox{arg} {\\mathcal F}_j]_{\\Gamma_{I+}}+ P\\,, \\qquad\n P = \\left\\{\n\\begin{array}{c}\n 2 \\quad \\mbox{if} \\quad \\mbox{tr}J>0 \\\\\n 0 \\quad \\mbox{if} \\quad \\mbox{tr}J<0 \n\\end{array}\n\\right. \\,.\n\\end{equation}\nBy using (\\ref{cyclic:hb_ex}) for the real and imaginary parts of \n${\\mathcal F}_j$, $[\\mbox{arg} {\\mathcal F}_j]_{\\Gamma_{I+}}$ is\neasily calculated numerically by a line sweep over\n$0<\\lambda_I<{\\mathcal R}$. Then, by using (\\ref{cyc:det_tr}) to\ncalculate $\\mbox{tr}(J)$, $P$ is readily determined. In this way, \n(\\ref{wind:last_2}) leads to a highly tractable numerical\nprocedure to calculate $N_j$. This criterion was used for all the\nresults below to identify regions in parameter space where\ninstabilities occur away from any Hopf bifurcation boundary.\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{fig_2spot_34.eps}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{fig_2spot_34_large.eps}\n\\caption{Hopf bifurcation boundaries in the $\\tau$ versus $D$ plane\n for $m=2$, $r_0=0.25$, and with parameters as in (\\ref{cyc:par}), \n computed from (\\ref{cyclic:hb_ex}). Left panel: the heavy solid\n curve and the solid curve are the Hopf bifurcation boundaries for\n the synchronous and asynchronous modes, respectively. Inside the\n respective lobes the corresponding mode is linearly unstable, as\n verified by the winding number criterion (\\ref{wind:last_2}). Right\n panel: same plot except that we include the Hopf bifurcation\n boundary for the synchronous mode from the leading-order\n $D={D_0\/\\nu}\\gg 1$ theory, computed from (\\ref{dlarge:sync_new}).}\n \\label{fig:m2}\n\\end{center}\n\\end{figure}\n\nIn Fig.~\\ref{fig:m2} we plot the Hopf bifurcation boundaries when\n$m=2$ and $r_0=0.25$. From the left panel of this figure, the\nsynchronous mode is unstable in the larger lobe shaped region, whereas\nthe asynchronous mode is unstable only in the small lobe for small\n$D$, which is contained within the instability lobe for the\nsynchronous mode. In the right panel of Fig.~\\ref{fig:m2} we show the\nHopf bifurcation boundary for the synchronous mode, as obtained from\n(\\ref{dlarge:sync_new}), corresponding to the leading-order\n$D={D_0\/\\nu}\\gg 1$ theory. Since the instability lobe occurs for only\nmoderate values of $D$, and $\\epsilon=0.05$ is only moderately small,\nthe leading-order theory from the $D={D_0\/\\nu}$ regime is, as\nexpected, not particularly accurate in determining the Hopf\nbifurcation boundary. The fact that we have stability at a fixed $D$\nfor $\\tau\\gg 1$, which corresponds to very fast intracellular\ndynamics, is expected since in this limit the intracellular dynamics\nbecomes decoupled from the bulk diffusion. Alternatively, if $\\tau\\ll\n1$, then for a fixed $D$, the intracellular reactions proceed too\nslowly to create any instability. Moreover, in contrast to the large\nregion of instability for the synchronous mode as seen in\nFig.~\\ref{fig:m2}, we observe that the lobe of instability for the\nasynchronous mode only occurs for small values of $D$, where the\ndiffusive coupling, and communication, between the two cells is rather\nweak. Somewhat more paradoxically, we also observe that the\nsynchronous lobe of instability is bounded in $D$. This issue is\ndiscussed in more detail below.\n\nIn Fig.~\\ref{fig:m2_all} we show the effect of changing the ring\nradius $r_0$ on the Hopf bifurcation boundaries. By varying $r_0$, we\neffectively are modulating the distance between the two cells. From\nthis figure we observe that as $r_0$ is decreased, the lobe of\ninstability for the asynchronous mode decreases, implying, rather\nintuitively, that at closer distances the two cells are better able to\nsynchronize their oscillations than when they are farther apart. We\nremark that results from the leading-order theory of \\S\n\\ref{sec:largeD} for the $D={\\mathcal O}(\\nu^{-1})$ regime would be\nindependent of $r_0$. We further observe from this figure that a\nsynchronous instability can be triggered from a more clustered spatial\narrangement of the cells inside the domain. In particular, for $D=5$\nand $\\tau=0.6$, we observe from Fig.~\\ref{fig:m2_all} that we are\noutside the lobe of instability for $r_0=0.5$, but inside the lobe of\ninstability for $r_0=0.25$ and $r_0=0.75$. We remark that due to the\nNeumann boundary conditions the cells on the ring with $r_0=0.75$ are\nclose to two image cells outside the unit disk, which leads to a\nqualitatively similar clustering effect of these near-boundary cells\nas when they are on the small ring of radius $r_0=0.25$.\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{fig_2spot_123456.eps}\n\\caption{Hopf bifurcation boundaries for the synchronous mode (larger\n lobes) and the asynchronous mode (smaller lobes) in the $\\tau$\n versus $D$ plane for $m=2$ and for three values of $r_0$, with\n $r_0=0.5$ (heavy solid curves), $r_0=0.75$ (solid curves), and\n $r_0=0.25$ (dashed curves). The other parameters are given in\n (\\ref{cyc:par}). We observe that as $r_0$ decreases, where the two\n cells become more closely spaced, the lobe of instability for the\n asynchronous mode decreases.}\n \\label{fig:m2_all}\n\\end{center}\n\\end{figure}\n\nIn Fig.~\\ref{fig:m3} we plot the Hopf bifurcation boundaries when\n$m=3$ and $r_0=0.5$. For $m=3$, we now observe that the region where\nthe synchronous mode is unstable is unbounded in $D$. The lobe of\ninstability for the asynchronous mode still exists only for small $D$,\nas shown in the right panel of Fig.~\\ref{fig:m3}. In this case, we\nobserve that the Hopf bifurcation boundary for the synchronous mode,\ncorresponding to the leading-order $D={D_0\/\\nu}\\gg 1$ theory and computed\nfrom (\\ref{dlarge:sync_new}), now agrees rather well with results\ncomputed from (\\ref{cyclic:hb_ex}).\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{fig_3spot.eps}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{fig_3spot_zoom.eps}\n\\caption{Left panel: Hopf bifurcation boundaries in the $\\tau$ versus\n $D$ plane for the synchronous mode for $m=3$ equally-spaced cells on\n a ring of radius $r_0=0.50$ (heavy solid curves), as computed from\n (\\ref{cyclic:hb_ex}), with parameters as in (\\ref{cyc:par}). The\n dashed curve is the Hopf bifurcation boundary from the leading-order\n $D={D_0\/\\nu}$ theory computed from (\\ref{dlarge:sync_new}). Right\n panel: The Hopf bifurcation boundaries for the asynchronous mode\n (solid curve) and the synchronous mode (heavy solid curve) shown in\n a magnified region of $D$. The asynchronous mode is linearly\n unstable only inside this small lobe, which lies within the\n unstable region for the synchronous mode.}\n \\label{fig:m3}\n\\end{center}\n\\end{figure}\n\n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{fig_hopf.eps}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{fig_large.eps}\n\\caption{Left panel: Hopf bifurcation boundaries in the $\\tau$ versus\n $D$ plane for the synchronous mode for $m=5$ equally-spaced cells on\n a ring of radius $r_0=0.25$ (solid curves) and radius $r_0=0.5$\n (heavy solid curves) concentric with the unit disk, as computed from\n (\\ref{cyclic:hb_ex}), with parameters (\\ref{cyc:par}). Right panel:\n Comparison of the Hopf bifurcation boundaries for the synchronous\n mode with $r_0=0.5$ (heavy solid curves), as computed from\n (\\ref{cyclic:hb_ex}), with that obtained from\n (\\ref{dlarge:sync_new}) for the leading-order $D={D_0\/\\nu}$ theory\n (solid curves). These curves agree well when $D$ is large.}\n \\label{fig:m5}\n\\end{center}\n\\end{figure}\n\nIn the left panel of Fig.~\\ref{fig:m5} we plot the Hopf bifurcation\nboundaries for the synchronous mode for $m=5$ when $r_0=0.5$ (heavy\nsolid curves) and for $r_0=0.25$ (solid curves). We observe that for\nmoderate values of $D$, the Hopf bifurcation values do depend\nsignificantly on the radius of the ring. The synchronous mode is\nunstable only in the infinite strip-like domain between these Hopf\nbifurcation boundaries. Therefore, only in some intermediate range of\n$\\tau$, representing the ratio of the rates of the intracellular\nreaction and bulk decay, is the synchronous mode unstable. As\nexpected, the two curves for different values of $r_0$ coalesce as $D$\nincreases, owing to the fact that the leading-order stability theory\nfor $D={D_0\/\\nu}\\gg 1$, as obtained from (\\ref{dlarge:sync_new}), is\nindependent of $r_0$. In the right panel of Fig.~\\ref{fig:m5} we\ncompare the Hopf bifurcation boundaries for the synchronous mode with\n$r_0=0.5$ with that obtained from (\\ref{dlarge:sync_new}),\ncorresponding to the leading-order theory in the $D={D_0\/\\nu\\gg 1}$\nregime. Rather curiously, we observe upon comparing the solid curves\nin the left and right panels in Fig.~\\ref{fig:m5} that the Hopf\nbifurcation boundaries from the $D={\\mathcal O}(1)$ theory when\n$r_0=0.25$, where the five cells are rather clustered near the origin,\nagree very closely with the leading order theory from the\n$D={D_0\/\\nu}\\gg 1$ regime. Since the clustering of cells is\neffectively equivalent to a system with a large diffusion coefficient,\nthis result above indicates, rather intuitively, that stability\nthresholds for a clustered spatial arrangement of cells will be more\nclosely approximated by results obtained from a large $D$\napproximation than for a non-clustered spatial arrangement of cells.\nIn Fig.~\\ref{fig:asy_m5} we plot the Hopf bifurcation boundaries for\nthe distinct asynchronous modes when $m=5$ for\n$r_0=0.5$ (left panel) and $r_0=0.75$ (right panel), as computed from\n(\\ref{cyclic:hb_ex}) with $j=2,5$ (larger lobe) and with $j=3,4$\n(smaller lobe). The asynchronous modes are only linearly unstable\nwithin these small lobes. \n\n\\begin{figure}[htbp]\n\\begin{center}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{fig_5asyn_r05.eps}\n\\includegraphics[width=0.45\\textwidth,height=5.0cm]{fig_5asyn_r075.eps}\n\\caption{Hopf bifurcation boundaries for the two distinct asynchronous\n modes when $m=5$ for $r_0=0.5$ (left panel) and $r_0=0.75$ (right\n panel), as computed from (\\ref{cyclic:hb_ex}) with $j=2,5$ (larger\n solid curve lobe) and with $j=3,4$ (smaller dashed curve lobe). The\n heavy solid curves are the Hopf bifurcation boundaries for the\n synchronous mode. The parameters are as in (\\ref{cyc:par}). The\n asynchronous mode for $j=2,5$ and $j=3,4$ is linearly unstable only\n inside the larger and smaller lobe, respectively.}\n\\label{fig:asy_m5}\n\\end{center}\n\\end{figure}\n\nTo theoretically explain the observation that the instability region\nin the $\\tau$ versus $D$ plane for the synchronous mode is bounded for\n$m=2$, but unbounded for $m\\geq 3$, we must first extend the large $D$\nanalysis of \\S \\ref{sec:odes} to the case of $m$ small cells. We\nreadily derive, assuming identical behavior in each of the $m$\ncells, that the reduced cell-bulk dynamics (\\ref{d:odes}) for one cell\nmust be replaced by\n\\begin{equation}\\label{d:modes}\n U_0^{\\prime}= -\\frac{1}{\\tau} \\left( 1 + \\frac{2\\pi m\n d_1}{|\\Omega|}\\right) U_0 + \\frac{2\\pi d_2 m}{\\tau |\\Omega|} u_1\n \\,, \\qquad \\vecb u^{\\prime} = \\vecb F(\\vecb u) + \\frac{2\\pi}{\n \\tau} \\left[d_1U_0-d_2u_1\\right] \\vecb e_1 \\,,\n\\end{equation}\nwhen there are $m$ cells. This indicates that the effective domain\narea is ${|\\Omega|\/m}={\\pi\/m}$ when there are $m$ cells. Therefore, to\nexamine the stability of the steady-state solution of (\\ref{d:modes})\nfor the Sel'kov model, we need only replace $|\\Omega|$ with ${|\\Omega|\/m}$\nin the Routh-Hurwitz criteria for the cubic (\\ref{largeD:3ode}).\n\nWith this approach, in Fig.~\\ref{bif:m3m5} we show that there are two\nHopf bifurcation values of $\\tau$ for the steady-state solution of\n(\\ref{d:modes}) when $m=3$ and $m=5$. These values correspond to the\nhorizontal asymptotes as $D\\to\\infty$ in Fig.~\\ref{fig:m3} for $m=3$\nand in Fig.~\\ref{fig:m5} for $m=5$. The numerical results from XPPAUT\n\\cite{xpp} in Fig.~\\ref{bif:m3m5} then reveal the existence of a\nstable periodic solution branch connecting these Hopf bifurcation\npoints for $m=3$ and $m=5$. A qualitatively similar picture holds for\nany $m\\geq 3$. In contrast, for $m=2$, we can verify numerically using\n(\\ref{largeD:3ode}), where we replace $|\\Omega|$ with ${|\\Omega|\/2}$,\nthat the Routh-Hurwitz stability criteria $p_1>0$, $p_3>0$, and\n$p_1p_2>p_3$ hold for all $\\tau>0$ when $m=2$ (and also $m=1$).\nTherefore, for $m=2$, there are no Hopf bifurcation points in $\\tau$\nfor the steady-state solution of (\\ref{d:modes}). This analysis\nsuggests why there is a bounded lobe of instability for the synchronous\nmode when $m=2$, as was shown in Fig.~\\ref{fig:m2}.\n\nWe now suggest a qualitative reason for our observation that the lobe\nof instability for the synchronous mode is bounded in $D$ only when\n$m0$~\\cite{PRL01}. In this limit the Hund and JT intra-molecular inter-orbital \ninteraction terms can be directly combined in the form (\\ref{H_el-ph}) \nwith $J= -J_H + J_{JT}$. In $C_{60}$, $J_H \\simeq 0.03-0.1~\\mathrm{eV}$~\\cite{Martin,Manini02}, \nwhereas $J_{JT} \\simeq 0.06-0.12~\\mathrm{eV}$~\\cite{Schluter,Varma,Gunnarsson1995,Manini94,Manini02}.\nThe total result -- and the only one compatible with s-wave superconductivity, \nwith a spin 1\/2 Mott insulator, and with a moderate spin gap near 0.1 eV --\nis a relatively weak unretarded attraction which has the form of an \n{\\em inverted} Hund's rule coupling~\\cite{Science,Granath}. This is the\napproximation we shall adopt, keeping in mind that it is strictly valid only when \n$ZW < \\hbar\\omega$, that is, close enough to the Mott transition. \n\nThe lattice expansion characterizing the expanded fullerides is believed to\nhave little effect on either $ J_H$ or $J_{JT}$. On the other hand, expansion will surely \ndecrease $W$ and increase $U$, so it can be modeled as a gradual increase of $U\/W$, \nreflecting both the band narrowing due to smaller overlap between molecular \nwave functions, and a decreased screening strength. \nAs discussed above, this Hamiltonian model, even if not really simple, has many body \ninteractions that are strictly on site, an ideal situation for Dynamical Mean-Field \nTheory (DMFT)~\\cite{DMFT}, one of the most popular and powerful tools in the field of \nstrongly correlated electron systems that we briefly describe in the next subsection.\n\n\\subsection{Dynamical Mean-Field Theory}\n\\label{Sec:dmft}\nDMFT is a quantum version of classical mean-field theory, which provides \nan exact description of the local dynamics, at the price of freezing away all spatial fluctuations.\nThe mean-field scheme is formulated by a mapping of the lattice model onto an Anderson \nimpurity model (AIM) embedded in a free-electron Fermi ``bath'' subject to a self-consistency condition~\\cite{DMFT}.\nIn our model, the effective AIM is three-fold\norbitally degenerate, with p-like levels representing the $t_{1u}$ orbitals,\neach hybridized with a bath. The Hamiltonian is \n\\begin{equation}\n\\label{aim}\n{\\cal H} = {\\cal H}_U +{\\cal H}_J +\\sum_{ka\\sigma} \\varepsilon_{ka} c^{\\dagger}_{ka\\sigma}c^{\\phantom{\\dagger}}_{ka\\sigma}+\\sum_{ka\\sigma} \nV_{ka} (c^{\\dagger}_{ka\\sigma}p^{\\phantom{\\dagger}}_{a\\sigma} + H.c.),\n\\end{equation}\nwhere ${\\cal H}_U$ and ${\\cal H}_J$ are (\\ref{H-U}) and (\\ref{H_el-ph}) for fermions on the impurity orbitals, \n$\\varepsilon_{ka}$ are the bath energy levels labeled by the index $k$ and by an orbital index $a=x,y,z$, \n$V_{ka}$ are the hybridization parameters between the bath fermions, created by $c^{\\dagger}_{ka\\sigma}$, \nand the impurity fermions, created by $p^{\\dagger}_{a\\sigma}$.\nThe mean-field scheme implies a self-consistency condition that depends on the \nimpurity Green's function and on the bare density of states of the original lattice. \nThroughout our calculations, we use an infinite-coordination Bethe lattice, \nwhose density of states is semicircular. This is a reasonable description of a \nrealistic three-dimensional density of states devoid of accidental features \nsuch as van-Hove singularities. It is moreover particularly convenient since it leads \nto a very simple and transparent form of the self-consistency condition.\nThe Bethe lattice is the $z\\to\\infty$ limit of a Cayley tree of coordination \n$z$, scaling the nearest-neighbor hopping in each one of the $z$ directions as \n$t\/\\sqrt{z}$. The resulting semicircular density of states has bandwidth $=4t$.\n \nThe self-consistency condition requires the so-called Weiss field \n(i.e., the non-interacting Green's function of the AIM)\n\\begin{equation}\n\\label{hybridization}\n{\\cal G}_0^{-1}(i\\omega_n)_a=(i\\omega_n+\\mu) - \\sum_k\\frac{V_{ka}^2}{i\\omega_n-\\varepsilon_{ka}},\n\\end{equation}\nto be related to the local interacting Green's function, $G_a(i\\omega_n)$, not only \nthrough the Dyson equation for the AIM\n\\[\n{\\cal G}_0^{-1}(i\\omega_n)_a = G_a^{-1}(i\\omega_n) + \\Sigma_a(i\\omega_n),\n\\]\nwhich requires the full solution of the impurity model, but also by the additional \nselfconsistency equation\n\\begin{equation}\n\\label{selfconsistency}\n\\mathcal{G}_{0a}^{-1}(i\\omega_n)=(i\\omega_n+\\mu) -t^2G(i\\omega_n)_a.\n\\end{equation}\n\nDMFT for a given model thus amounts to solve iteratively the AIM until the impurity Green's function satisfies \nEq.~(\\ref{selfconsistency}). In this work, we solved the threefold degenerate AIM by exact diagonalization.\nThat requires truncating the sum over $k$ in Eqs.~(\\ref{aim}) and (\\ref{selfconsistency}) to a finite \nand relatively small number of baths $N_b$, so that the Hamiltonian can be \ndiagonalized in the finite resulting Hilbert space. We generally used $N_b=4$ for each orbital.\nMost of the results we shall present are at zero temperature, where we can use the \nLanczos algorithm to calculate the Green's function without fully diagonalizing the Hamiltonian. \nThe finite temperature results for the specific heat and its jump at $T_c$ reported in \nSec.~\\ref{Sec:Comparison} are the exception. They are obtained by means of the \nfinite-temperature extension of Lanczos~\\cite{caponelanczos},\nwhere the thermal Green's function is expressed as a sum over the low-lying eigenvectors \n$\\vert n\\rangle$ and eigenvalues $E_n$ of the impurity model\n\\begin{equation}\\label{G_sum_Gm}\nG_{a\\sigma}(i\\omega_n)=\\frac{1}{Z}\\sum_m e^{-\\beta E_m} G_{a\\sigma}^{(m)}(i\\omega_n)\n\\end{equation}\nwith\n\\begin{equation}\\label{G_m}\nG_{a\\sigma}^{(m)}(i\\omega_n)\\equiv \\sum_n\n\\frac{\\left\\vert \\langle n\\vert p_{a\\sigma}\\vert m\\rangle\\right\\vert ^2}{E_m-E_n-i\\omega_n} +\n\\sum_n \\frac{\\left\\vert \\langle n\\vert p^{\\dagger}_{a\\sigma}\\vert m\\rangle\\right\\vert ^2}\n{E_n-E_m-i\\omega_n}\\,.\n\\end{equation}\nThe Boltzmann factor in Eq. (\\ref{G_sum_Gm}) guarantees that the lower the temperature, \nthe smaller the number of excited states \nthat actually contribute to the Green's function. Therefore, for sufficiently low temperature, \nwe can still use the Lanczos algorithm to find the lowest energy eigenstates.\nIn practice, only a limited number of states, $20\\div 25$, can be calculated in \na reasonable time. Unfortunately, this number is insufficient to \nmake the truncation error in Eq.~(\\ref{G_sum_Gm}) negligible for the present model \nin the relevant correlated regime. We estimate the systematic error on the \nGreen's function to be of order of a few per cent, a level of accuracy that does \nnot allow to fully determine thermodynamic properties such as the critical temperature. \nLuckily though, this error affects much more the absolute specific heat value than \nits relative changes, including the superconducting jump in units of $T_c$ \nwhich we shall discuss further below (see Sec.~\\ref{Sec:Comparison}). We underline that \nthis limitation does not affect by any means the $T=0$ calculations.\n\nThe above DMFT equations refer to a paramagnetic phase where no symmetry breaking occurs. \nHowever, in this work we shall be crucially interested in $s$-wave superconducting \nand in antiferromagnetic solutions, where symmetry is broken.\nSuperconductivity is conveniently studied in the Nambu-Gor'kov representation \nby introducing spinors\n\\[\n\\psi_{\\mathbf{k} a} = \n\\left(\\begin{array}{c}\np^{\\phantom{\\dagger}}_{\\mathbf{k} a\\uparrow}\\\\\np^\\dagger_{-\\mathbf{k} a\\downarrow}\n\\end{array}\\right),\n\\]\n\nand defining accordingly the Green's function $G_{\\mathbf{k} a}(\\tau) = -\\langle T_\\tau\n\\left(\\psi^{\\phantom{\\dagger}}_{\\mathbf{k} a}(\\tau)\\,\\psi^\\dagger_{\\mathbf{k} a}(0)\\right)\\rangle$ in imaginary time as a $2\\times 2$ matrix \nthat satisfies the Dyson equation in Matsubara frequencies\n\n\\begin{equation}\n\\label{dyson}\nG_{\\mathbf{k} a}(i\\omega) = G^0_{\\mathbf{k} a}(i\\omega)+ G^0_{\\mathbf{k} a}(i\\omega)\\,\\Sigma(i\\omega)\\,G_{\\mathbf{k} a}(i\\omega),\n\\end{equation}\nwith $G^0_{\\mathbf{k} a}(i\\omega)$ the non-interacting value. The single-particle self-energy $\\Sigma(i\\omega)$ \nis also a $2\\times 2$ matrix whose off-diagonal element $\\Delta(i\\omega)$, when finite, \nsignals a superconducting phase. The DMFT self-consistency can be written now as\n\\begin{equation} \n\\label{selfsuper} \n\\mathcal{G}_0^{-1}(i \\omega_n)= i\\omega_n \\,\\tau_0 + \\mu \\,\\tau_3 \n -t^2 \\, \\tau_3 \\,G(i\\omega_n)\\, \\tau_3,\n\\end{equation} \nwhere $\\tau_{0}$ and $\\tau_{3}$ are Pauli matrices. \n\nAnalogously, antiferromagnetism in a bipartite lattice is conveniently described using the spinor \n\\[\n\\psi_{\\mathbf{k} a \\sigma} = \n\\left(\\begin{array}{c}\np^{\\phantom{\\dagger}}_{\\mathbf{k} a\\sigma}\\\\\np^{\\phantom{\\dagger}}_{\\mathbf{k} + \\mathbf{Q} a\\sigma}\\end{array}\\right),\n\\]\nwith $\\mathbf{k}$ in the magnetic Brillouin zone and $\\mathbf{Q}$ the modulation vector. \nThis leads once more to a $2\\times 2$ Green's function and self-energy matrices related by the Dyson equation \n(\\ref{dyson}). Here too, a finite off-diagonal element signals an antiferromagnetically ordered phase.\nThe self-consistency equation for antiferromagnetism exploits the bipartite property of the lattice. \nIndicating one sublattice with A and the other with B, the general self-consistency condition is \n\\begin{equation}\n\\label{selfconsistencyafm}\n{\\cal G}_0^{-1}(i\\omega_n)_{\\sigma A}=(i\\omega_n+\\mu)-t^2\\,G(i\\omega_n)_{\\sigma B},\n\\end{equation}\nwhich becomes Eq.~(\\ref{selfconsistency}) if $G_{\\sigma A} \\equiv G_{\\sigma B}$, i.e. if \nthe system is nonmagnetic. When the system is antiferromagnetic, then $G_{\\sigma A}\\equiv G_{-\\sigma B}$. \nThus we can eliminate the sublattice B from Eq. (\\ref{selfconsistencyafm}), and obtain the following \nresult for the self-consistency equation: \n\\begin{equation}\n\\label{selfconsistencyafm-bis}\n{\\cal G}_0^{-1}(i\\omega_n)_{\\sigma A}=(i\\omega_n+\\mu)-t^2G(i\\omega_n)_{-\\sigma A}.\n\\end{equation}\nThis equation is valid for a Bethe lattice with nearest neighbor hopping, a case with \nperfect nesting that is rather exceptional in realistic antiferromagnets. \nA way to simulate imperfect nesting typical of more realistic situations, while still taking \nadvantage of the Bethe-lattice simplifications, is to add a next-nearest-neighbor \nhopping $t^\\prime\/z$ in the Cayley tree (merely a device to eliminate nesting, \nnot meant to suggest next nearest neighbor hopping, small in fullerides). In the limit \n$z\\to\\infty$ of the Bethe lattice, the self-consistency equation becomes\n\\begin{eqnarray}\n{\\cal G}_0^{-1}(i\\omega_n)_{\\sigma A}&=& (i\\omega_n+\\mu)-t^2\\,G(i\\omega_n)_{\\sigma B}-{t^{\\prime}}^2\\,G(i\\omega_n)_{\\sigma A}\n\\nonumber\\\\\n&=& (i\\omega_n+\\mu)-t^2\\,G(i\\omega_n)_{-\\sigma A}-{t^{\\prime}}^2\\,G(i\\omega_n)_{\\sigma A}.\n\\label{selfconsistency2}\n\\end{eqnarray}\n\nFor both broken-symmetry phases, if the diagonal elements of the self-energy matrix \nat low Matsubara frequencies follow the conventional Fermi-liquid behavior -- which \nwe always find to be the case, \n\\begin{equation}\n\\Sigma_{diagonal}(i\\omega_n) \\simeq \\left(1-\\frac{1}{Z}\\right)\\,i\\omega_n,\n\\label{zeta}\n\\end{equation}\nthe actual value of the spectral gap in the single particle spectrum is given by \n$Z\\,\\Delta(0)$, the zero frequency anomalous (superconducting or antiferromagnetic) \nself-energy multiplied by the so-called quasiparticle weight $Z$.\n\nWe end by noting that, within DMFT, one can search for solutions with different \nsymmetries by simply allowing\/preventing symmetry breaking order parameters, even \nin regions where the chosen phase is not the most stable. When two or more \nsolutions coexist, the stable one is determined by an explicit energy \ncalculation. As we shall discuss, for a wide range of $U\/W$ values \nwe do find coexisting superconducting and antiferromagnetic solutions,\nthe former prevailing at smaller $U$, the latter at larger $U$. \nThe physical phase diagram thus exhibits a first-order transition \nbetween these two symmetry broken phases -- a nonmagnetic $s$-wave superconductor,\nand an insulating spin 1\/2 antiferromagnet -- taking place when \nthe respective energy curves intersect.\n\n\n\\subsection{$T =$ 0 Phase Diagram}\nModeling lattice expansion of fullerides as a gradual increase of $U\/W$, \nwe can proceed to analyze the theoretical zero temperature ``phase diagram'' of our model\nobtained by DMFT as a function of $U\/W$.\nStarting with the uncorrelated system with $U = 0$ the model initially exhibits \nstraight BCS superconductivity driven by JT phonons \n(i.e., the attractive $J$ we introduced above) with an $s$-wave ($S = L = 0$) \norder parameter\\footnote{Since gauge symmetry is broken, we are allowed to assume $P_{SC}$ real} \n\\begin{equation}\nP_{SC} = \\frac{1}{N}\\sum_i\\sum_{a=x,y,z}\\,\n\\langle p^\\dagger_{i\\,a\\uparrow}\\, p^\\dagger_{i\\,a\\downarrow}\\rangle,\n\\label{Cooper-channel}\n\\end{equation}\nwhere $N$ is the number of molecules and $p^\\dagger_{i\\,a\\sigma}$ creates an electron on molecule $i$, with spin $\\sigma$ and in orbital \n$a=x,y,z$. The Fermi-liquid scattering amplitude in \nthe Cooper channel, measuring the strength of\nthe effective attraction, is $A=-10 J\/3$. We note that owing to fairly strong JT interactions\n~\\cite{Manini94, Gunnarsson1995}, if Hund's exchange $J_H$ were neglected, the\ndimensionless JT coupling constant of fullerene controlling superconductivity would be \nnumerically very large, $\\lambda = \\rho_0\\vert\\,A\\, \\vert\\simeq 0.6 - 1.0$, where \n$\\rho_0$ ($\\simeq 2.4~\\mbox{eV}^{-1}$) is the bare density of states per spin and orbital.\nTurning on a weak Coulomb repulsion $U$ on top of that will reduce the pairing attraction \nin this regime. Perturbatively one obtains for small $U$ that $A= -10\\, J\/3 + U$. \nSince $J$ is insensitive to expansion while $U\/W$ increases, this implies \nthat in this picture, where Hund's rule exchange is neglected, $T_c$ should \nalways {\\it decrease} upon expansion, a prediction which is at odds with experiments. \n\nIn fact, as anticipated, Hund's rule exchange is not negligible, and its \neffect is to introduce a substantial cancellation in $J$ leading to a largely\nreduced effective coupling $\\lambda_{eff} \\simeq \\frac{10}{3}\\left(J_{JT}-J_H \\right)\\,\\rho_0 $.\nShould we interpret the ubiquitous spin gap of 0.1~eV~\\cite{Zimmer,Brouet,Brouet-1} as \nthe C$_{60}^{n-}$ molecular excitation energy between its low spin ground state \nwith $S= 1\/2$ and $S=0$ for $n=3$ and $n=4$, respectively, and its high spin excited \nstate, with $S= 3\/2$ and $S=1$ for $n=3$ and $n=4$, respectively, we would conclude \nthat, to be consistent with our model where both gaps are equal to $5J$~\\cite{Capone04}, \nthe total effective inverted exchange $J$, comprehensive of both JT and Hund's rule\nexchange, should be approximately $J=0.02$~eV. In reality, the \nqualitative scenario we will describe is relatively independent of the precise \nvalue of $J$ provided it is inverted, i.e., negative and small. \n\nOnce exchange is included, the resulting $\\lambda_{eff} \\simeq 0.16$ is now much \nsmaller -- in fact way too small to explain within conventional BCS or Migdal-Eliashberg theory \nany of the observed values of $T_c$ in fullerides (let alone the non-monotonic behavior \nof $T_c$ versus the density of states for expanded fullerides~\\cite{Dahlke,Ganin}). \nThings get worse when we increase the on-site Coulomb repulsion $U$ closer to\nrealistic values, $U\\sim W$ and beyond. In the conventional weakly correlated picture \n$U$ would provide in the electron pairing problem a repulsive ``Coulomb pseudopotential'' whose \nbare value is $\\mu_* = U\\,\\rho_0\\simeq 3$~\\cite{BCSreview}. Simply comparing \nthese bare values of $\\lambda$ and $\\mu_*$ we should conclude that $s$-wave BCS \nsuperconductivity in fullerenes is simply impossible (with the obvious proviso that for small $U$ \nan unretarded treatment of JT phonon interactions is not really justified).\n\nThe full DMFT solution of the model for $J=0.05~W$ and increasing $U\/W$ \nyields the phase diagram in Fig.~\\ref{phase-diagram}. While confirming the above\nexpectations for moderate $U$, it has a surprise in reserve at \nlarger $U$ values, where the Mott transition is approached.\n\n\\begin{figure}\n\t\\includegraphics[width=8.5cm]{fig2.eps}\n\t\\caption{\\label{phase-diagram} Superconducting solution (triangles) and antiferromagnetic solution (squares) \n\tzero-frequency anomalous self-energies $\\Delta(0)$ as function of $U\/W$. Solid symbols are used when the \n corresponding symmetry broken phase is stable, while open symbols when it is metastable, i.e. has lower energy \n then the other phase. The first-order transition between the two phases is indicated by a \n vertical line separating the superconductor from the antiferromagnetic Mott insulator. \n Panel (b) corresponds to the case in the presence of a frustrating next-nearest neighbor \n hopping $t'=0.3 t$, absent in (a).}\n\t\\end{figure} \n\nFig.~\\ref{opdelta} shows the zero frequency anomalous single-particle self-energies calculated within DMFT \nfor the superconducting and for the antiferromagnetic solution. At $U=0$, the model is \nan $s$-wave BCS superconductor, with an exponentially small superconducting $\\Delta(0)$. \nIt is too small to be visible in the figure, since the effective exchange-reduced $\\lambda \\sim 0.2$ \nis as was said very weak. Beginning from zero, the increase of $U$ first rapidly \ndestroys the weak BCS superconductivity. The superconducting $\\Delta(0)$ vanishes at \nroughly the mean-field value $U=10\/3 J$, and above this value of $U$ the ground state \nbecomes a normal metal as expected. Upon further increasing $U\/W$ the model remains \na normal metal -- no superconductivity, no antiferromagnetism. However, the importance \nof electron correlations increases with $U$, as signaled for example in the DMFT \nspectral function (not shown) by the gradual\nformation of incoherent Hubbard bands on both sides of the Fermi level. The metallic \ncharacter persists until a hypothetically continuous Mott transition eventually \nreached near $U\/W \\sim 1.5$, where $Z=0$ and the metallic character is extinguished. \n\nBefore this happens however, $s$-wave superconductivity re-enters from the normal metal\nstate. The anomalous self energy $\\Delta(0)$, proportional to the superconducting \norder parameter has, as a function of $U$, a bell-shaped behavior --\na ``superconducting dome'' as it is called in cuprates -- hitting a large maximum before dropping again. The \nre-entrant superconductive behavior is a clear realization of phonon-induced strongly \ncorrelated superconductivity (SCS)~\\cite{Science}. The sharply rising order parameter\nedge with increasing $U\/W$ can in our view explain the strong rise of $T_c$ upon \nlattice expansion in non-expanded compounds, previously (and we believe incorrectly)\nattributed to a BCS-like increase of density of states upon band narrowing. \nPast the dome maximum, and upon increasing expansion, the SCS superconducting order\nparameter declines, and would eventually drop to zero at the continuous metal \ninsulator transition near $U\/W \\sim 1.5$. This continuous decline of superconductivity\nis preempted by a first order transition to a lower energy antiferromagnetic \nMott insulating phase, with order parameter\n\\begin{equation}\nP_{AFM}=\\frac{1}{N}\\sum_{i}(-1)^i(n_{i\\uparrow}-n_{i\\downarrow}),\n\\end{equation}\nwhere $n_{i\\sigma}=\\sum_{a=x,y,z} p^\\dagger_{i\\,a\\sigma}p^{\\phantom{\\dagger}}_{i\\,a\\sigma}$ is the full occupation number \nwith spin $\\sigma$ at molecule $i$. \nThe exact location of the superconductor-insulator transition depends on details. \nFor strong nesting ($t^\\prime$=0) it takes place even before the superconducting dome maximum. \nIn Fig.~\\ref{opdelta} we show $\\Delta(0)$ of Fig.~\\ref{phase-diagram} \nin comparison with the spectral gaps $Z\\Delta(0)$, as well as the order \nparameters $P_{SC}$ and $P_{AFM}$. Notice that $Z$ for the superconductor is smaller and vanishes at the continuous \nmetal insulator transition, while $Z$ for the antiferromagnet is of order 1. In both \ncases the dimensionless order parameter essentially follows the behavior of the spectral gap.\n\n\\begin{figure}\n\\includegraphics[width=8.5cm]{fig3.eps}\n\\caption{\\label{opdelta} Superconducting solution (SC, left) and antiferromagnetic solution (AFM, right) \nanomalous self-energies (top) and order parameters (bottom) as a function of $U\/W$. \nThe top panels also show (green diamonds) the spectral gaps obtained multiplying $\\Delta$ by the quasiparticle \nweight $Z$ of each solution. The gaps are in units of the bandwidth $W$ (the order parameters are \nby definition dimensionless). Dashed vertical lines mark the first-order \nphase transition between the two solutions.}\n\\end{figure} \n\nWhen we add a next-nearest neighbor hopping $t^\\prime$ to mimic imperfect nesting \n(as we expect to find generically for realistic band structures, \nin particular in the face-centered cubic A$_3$C$_{60}$ materials), \nwe find that the superconducting phase is only weakly affected but antiferromagnetism is \nstrongly frustrated. As a result the superconducting region expands at expense of the \nmagnetic insulator, and the superconducting dome may emerge in full \n[panel (b) of Fig.~\\ref{phase-diagram}]. We propose that the gradual drop of superconducting\norder parameter past the dome maximum now naturally explains the decline of $T_c$\nof expanded fullerides~\\cite{Dahlke&Rosseinsky, Durand2003,Ganin}.\n\nFinally, past the first-order Mott transition, we find that the antiferromagnetic insulator \nis formed mainly by spin-1\/2 local configurations, which is in agreement with experiments \nin NH$_3$K$_3$C$_{60}$~\\cite{Iwasa}. We also predict that ambient pressure \nA15 Cs$_3$C$_{60}$, yet to be characterized, should similarly be a spin-1\/2\nantiferromagnetic insulator. Besides spin rotational symmetry, this kind of\nstate also breaks orbital rotational symmetry, signaling that spin ordering \nmust be generally accompanied by orbital ordering. In ammoniated fullerides that\nagain is consistent with experiment~\\cite{Iwasa}. We conclude that, in spite of strong \nsimplifying assumptions, our model seems able to reproduce very important features of the \nphase diagram of expanded fullerides. In the following we shall discuss in more\ndetail the strongly correlated superconducting phase near the Mott transition, and \nalso propose experiments that might distinguish it from a standard BCS state.\n\n\n\\section{Understanding Strongly Correlated Superconductivity from DMFT and the Impurity Model}\nThe re-emergence of phonon-driven superconductivity close to the Mott transition -- \nStrongly Correlated Superconductivity (SCS) -- was discussed in \\textcite{Science} \nin terms of Fermi-liquid theory. A key point of that phenomenon is the renormalization of \nthe effective bandwidth and thus of the effective mass, both controlled (in a Bethe \nlattice) by the quasiparticle weight $Z$, Eq.~(\\ref{zeta}). $Z(U)$ decreases as a function of\n$U\/W$ and vanishes at the continuous metal insulator transition point $U=U_c$, where the effective \nmass $m^*\/m = 1\/Z(U)$ diverges, and the effective quasiparticle bandwidth $W_* = ZW$ vanishes.\nAn estimate of the interaction between charged quasiparticles requires the evolution of fluctuations that \ntake place in charge space. Because charge fluctuations are gradually frozen away near the Mott transition, \nthe effective repulsion between quasiparticles is also renormalized down to some smaller value $U_* T_K$ the Kondo screening is lost, and that \nwas shown to imply a non Fermi-liquid phase characterized by a pseudogap in the \nsingle-particle spectrum and by several other singular properties~\\cite{DeleoPRL}. \nA very qualitative description of this phase is that, unlike the Kondo phase, \ntwo spins out of three pair off antiferromagnetically at any given time, leaving \nout a single spin 1\/2 available for Kondo screening. However, since orbital degeneracy is unbroken, \nthis residual spin $S=1\/2$ also carries orbital momentum $L=1$, which corresponds to an overscreened \nnon-Fermi liquid situation~\\cite{DeleoPRL}. \nIn this regime, it was predicted that the impurity contributions\nto the specific-heat coefficient and to the pair susceptibility in the $s$-wave \nchannel (\\ref{Cooper-channel}) diverge as $T^{-1\/5}$ at low temperature $T$. \nIn addition, the conduction electron scattering rate has a non-analytic \ntemperature-behavior $T^{2\/5}$. The local response functions to either\na quadrupolar field, which splits the orbital degeneracy, or to a magnetic \nfield which polarizes both spins and orbitals also diverge as $T^{-1\/5}$. \nThe two phases, the Kondo screened phase and the pseudogap phase, are separated \nby a critical point at $J = J_{cp} \\simeq T_K$. It is endowed with a finite entropy \n$1\/2\\,\\ln 3$, and with a divergent superconducting susceptibility with an exponent $1\/3$. \nAs shown in Fig.~\\ref{DOS}, the single-particle spectral function displays strong \ndeviations from a normal metal in the pseudogap phase and at the critical point.\n\n\\begin{figure}[t]\n\\includegraphics[width=7cm]{fig6.eps}\n\\caption{\\label{DOS} Impurity spectral function with \na rigid bath characterized by a flat density of states (no DMFT self-consistency) across \nthe critical point: the curves from top to bottom \ncorrespond to the evolution from the Kondo screened regime towards the pseudogap regime. \nThe dotted line is the spectral function \ndeep inside the non-Fermi-liquid phase $J\\gg T_K$. The top inset shows the same curves plotted on a \nlarger energy range, where the pseudogap of the dotted line is visible. The Fermi liquid\nbehavior corresponds to the normal state of the SCS superconducting phase for $J < J_{cp}$ \n(analogous to the overdoped regime of cuprates) whereas the pseudogap behavior \ncorresponds to the SCS normal state for $J > J_{cp}$ (analogous to the underdoped cuprates). \nFigure from \\textcite{DeleoPRL}.}\n\\end{figure}\n\nNear this critical point it has been found that the low-energy dynamics around the impurity is controlled \nby two separate energy scales~\\cite{DeleoPRL}, $T_+$ and $T_-$, \nwhose behavior is very different as a function of increasing $U\/W$. \nA higher energy scale $T_+$ is set by the critical $J_{cp}\\sim T_K$ and represents \nthe width of a broad incoherent resonance; it evolves smoothly and uneventfully\nas the critical point is crossed (See Fig.~\\ref{DOS}). A lower energy scale \n$T_-\\propto |J-T_K|^3$ measures instead the distance from the critical point, \nand leads simultaneously to a narrow resonance in the Fermi-liquid region, and to \nan equally narrow spectral density dip (the ``pseudogap'') in the pseudogap \nphase~\\cite{DeleoPRL}. Since in this phase the impurity still carries a residual \nspin-1\/2, there remains a finite value of the spectral function at the chemical \npotential, and the gap is not complete. The pseudogap widens if $J$ is increased, \nand the cusp-like dip in the impurity spectral function smoothly turns into a \ncusp-like peak, the value at the chemical potential staying fixed and constant. \nThis behavior is shown by the dotted curve in Fig.~\\ref{DOS} corresponding to a \nvery large pseudogap (see top inset), possessing a very tiny peak at the \nchemical potential. This indicates the existence of yet another energy scale \nbesides $T_+$ and $T_-$ that sets the width of the cusp peak. \n\n\n\\subsection{Anderson Impurity with a Self-Consistent Bath in DMFT}\n\nThe rigid bath AIM behavior and its critical point briefly reviewed above \nprovide a guide to the DMFT results once the impurity-bath coupling is self-consistently \ndetermined. First of all, since $T_K$ coincides within DMFT with the renormalized $ZW$ \nwhich in turn vanishes when the continuous Mott transition is approached, \nthe impurity critical point is inevitably met before the metal insulator transition\nas $U\/W$ is increased, at some $U_{cp} \\alt U_c$. This entails several important consequences:\n\n\\begin{itemize}\n\\item The normal state may be a Fermi liquid only far below the continuous metal insulator\ntransition, such as perhaps may be the case in the non-expanded fullerides. Expanded \ncompounds on the other hand are expected to have a non Fermi-liquid normal state, \nand eventually a pseudogap, possibly developing before the first order transition \nto the antiferromagnetic insulator.\n\n\\item The SCS superconducting pocket near the Mott transition reflects the leading \ninstability of the impurity critical point. In other words SCS superconductivity \nis the way in which the lattice model responds to impurity criticality and avoids it.\n\n\\item The lower energy scale vanishes right at the critical point, $T_-=0$. Here \n$T_+\\simeq J$ thus remains as the only energy scale controlling the magnitude of the \nsuperconducting energy gap. Away from $U=U_{cp}$, $T_-\\not =0$, and the amplitude of the \nsuperconducting gap should decrease monotonically with $(T_+-T_-)\/T_+$~\\cite{Capone04}, since \n$T_-$ cut-offs the local pairing instability. Therefore the gap should be maximum right at the impurity critical point (the top of the dome).\n\n\\item Even though the normal phase is non-Fermi liquid, well defined Bogoliubov \nquasiparticles should exist inside the SCS superconducting pocket. \n\n\\end{itemize}\nThe last statement comes from the fact that, at the impurity critical point, superconductivity \nprovides a new screening channel which helps the system get rid of the finite \nresidual entropy at the critical point, thus eliminating non-Fermi liquid singularities~\\cite{DeleoPRL,SchiroPRB}. \n\n\\begin{figure}[t]\n\\includegraphics[width=7cm]{fig7.eps}\n\\caption{\\label{Drude} Calculated weights of the zero-frequency delta-function contribution \nto the optical conductivity for the superconducting phase (superfluid stiffness, squares) \nand the normal phase (Drude weight, circles). The zero-frequency anomalous self-energy is also plotted.\nNote the higher values in the SCS superconductor, relative to the nonsuperconducting metal phase.}\n\\end{figure}\n\nThe role of superconductivity as a novel screening channel close to the \nimpurity model critical point should reflect, in the lattice model, into \na gain of band energy (the tight binding ``kinetic energy'') at the onset of superconductivity. \nOwing to a sum rule connecting kinetic energy and zero-frequency optical\nconductivity (sometimes referred to as ``Drude weight'')~\\cite{Drude}\nthe onset of SCS superconductivity close to the Mott transition\nshould lead to a Drude weight increase. This prediction is \nwell borne out by the full DMFT solution of our Hamiltonian.\nIn Fig.~\\ref{Drude} we plot the $\\omega=0$ (d.c.) optical conductivity of \nour model superconductor (where it coincides with the superfluid stiffness), defined by~\\cite{toschi}\n\\begin{equation}\nD_s= -E_{kin} + \\chi_{jj}({\\bf q} \\rightarrow 0, \\Omega=0),\n\\label{ds}\n\\end{equation} \nwhere $E_{kin}$ is the kinetic energy and $\\chi_{jj}$ the static limit of the paramagnetic\npart of the electromagnetic kernel\n\\begin{eqnarray}\n\\chi_{jj} & = & \\frac{2}{\\beta} \\sum_{n} \\int \\, d\\epsilon \\, N(\\epsilon) \\, \nV(\\epsilon) \\nonumber\\\\\n&\\times & \\left[ G(\\epsilon,\\omega_n) G^{*}(\\epsilon, \\omega_n) +\nF(\\epsilon,\\omega_n)F(\\epsilon,\\omega_n) \\right], \n\\label{chijj} \n\\end{eqnarray} \nwhere $V(\\epsilon)=(4t^2-\\epsilon^2)\/3$ is the current vertex in the Bethe lattice, while $G(\\epsilon,\\omega_n)$ and \n$F(\\epsilon, \\omega_n)$ are the normal and anomalous lattice Green functions, respectively. \nIn the same figure we also plot the zero-frequency d.c. conductivity (the Drude weight) of the \nunderlying metastable solution where superconductivity is \ninhibited, this state meant to provide a cartoon of the real normal phase above $T_c$. \nThe Drude weight is given by Eq. (\\ref{ds}) and \n(\\ref{chijj}) with $F(\\epsilon,\\omega_n)\\equiv0$. Upon entering the SCS dome \nfrom the low $U\/W$ side, the superconducting phase initially loses kinetic energy \nover the normal state as in ordinary BCS theory. However, upon increasing $U\/W$ at \nand beyond the dome maximum, the loss reverts over to a gain, and indeed \nmost of the SCS superconducting pocket is predicted to have a larger \nweight of the zero-frequency optical absorption than the non superconducting state.\nThe same behavior is displayed (Fig.~\\ref{kinetic.vs.potential}) by the energy balance between \nthe two solutions. Only far below the Mott transition the superconductor is stabilized \nby a potential energy gain, as is the case in BCS theory. In the pseudogap regime,\ncorresponding to the expanded fullerides near the Mott transition, the stabilization \nis associated with a kinetic energy gain.\n\\begin{figure}[t]\n\\includegraphics[width=7cm]{fig8.eps}\n\\caption{\\label{kinetic.vs.potential} Energetic balance underlying superconductivity. $\\Delta E_{pot} = E^{S}_{pot} - E^{N}_{pot}$ \nis the difference between the potential energies of the superconducting and the normal solution, \nwhile $\\Delta E_{kin} = E^{S}_{kin} - E^{N}_{kin}$ is the same difference between the kinetic energies of the two solutions.}\n\\end{figure}\n\nA similar phenomenon is well known in the optical conductivity of \nhigh-T$_c$ copper oxides~\\cite{VanDerMarel}. Our calculations show that an \nincrease in the zero-frequency optical conductivity or a kinetic \nenergy gain do not actually exclude an electron-phonon pairing mechanism, \nbut rather demonstrates the key importance of strong electronic correlations. \nThanks to pairing, the motion of carriers in the superconducting phase is \nfacilitated with respect to the pseudogap nonsuperconducting metal. In that\nanomalous metal -- a non Fermi liquid-- the interaction constraints jam the \nfree propagation of quasiparticles, causing a kinetic energy cost, partly released\nwith the onset of superconductivity. It would be of extreme interest \nif the optical conductivity increase demonstrated in cuprates could be investigated\nin fullerides, both regular and expanded, since that would help \ndiscriminate between conventional BCS and SCS.\n\n\\begin{figure}[b]\n\\includegraphics[width=8cm]{fig9.eps}\n\\caption{\\label{specific-heat-jump} Specific heat jump, in units of \n$T_c$, $\\gamma_0$, as function of $U\/W$. For comparison we also show the behavior \nof $1\/Z(U)$, which should correspond to the same quantity if Fermi liquid theory were \nvalid.}\n\\end{figure}\n\n\n\\section{Discussion of Experimental Probes}\n\\label{Sec:Comparison}\nIn the following we discuss whether and how the dominance of strong correlation \nwe propose for the superconducting fullerides can be reconciled with the list of experimental facts \ngiven earlier, apparently supporting a conventional BCS behavior. The first quantity \nwe discuss is the specific heat jump at $T_c$. Within BCS theory, the specific \nheat jump $\\Delta C_V$ at $T_c$ is approximately \n\\begin{equation}\n\\frac{\\displaystyle \\Delta C_V }{\\displaystyle T_c} \\simeq 1.52\\,\\gamma_*,\n\\label{heat-jump}\n\\end{equation}\nwhere $\\gamma_*$ is the specific heat coefficient of the metallic phase ($C_V = \\gamma_* T$)), \nproportional to the mass enhancement $m^*\/m$, in our case $1\/Z$. \nThe measured specific heat jump leads, through (\\ref{heat-jump}), to an estimate \nof $\\gamma_*\\simeq 3\\gamma_0$. This is indeed a rather low and regular value which has been \ncommonly advocated as evidence of weak correlations in fullerenes. However, the standard \nargument only holds provided the normal phase is Fermi-liquid, which is not applicable \nclose to a Mott transition, independently of any models. \n\nIn Fig.~\\ref{specific-heat-jump} we plot the jump in $C_V$ at $T_c$ for \nincreasing $U\/W$, compared with the Fermi-liquid estimate $1\/Z$.\\footnote{As mentioned in Sec. \\ref{Sec:dmft}, \nfor the present three-orbital model it has not been practical \nto include in Eq. (\\ref{G_sum_Gm}) a number of excited states sufficient to make the truncation error \nin the Green's function negligible. \nFor the data reported in Fig.~\\ref{specific-heat-jump} we have an average error of 3-4\\% \nwhich does not allow us to determine $T_c$ with sufficient accuracy. Yet, the jump of the specific heat relative to $T_c$ turns out to be almost \nindependent on the truncation error and on the details of the calculations.}\nAfter a region \nwhere the calculated quantity closely follows $1\/Z$, $\\Delta C_V\/T_c$ flattens out \nand stays roughly constant around $4\\,\\gamma_0$ up to the Mott \ntransition, despite a diverging $1\/Z$ (See Fig.~\\ref{specific-heat-jump}). \nThis shows that the energy scale that controls the superconducting instability \nis constant near the Mott transition, consistent with the single-impurity prediction \nthat this scale should be identified by $T_+$, a quantity of order $J$. \nThe bottom line conclusion here is that a normally sized specific-heat jump does not imply\nBCS superconductivity in fullerides.\n\nAnother physical quantity which seemingly pointed towards weak \ncorrelations in fullerides is the magnetic susceptibility measured in the normal \nphase, apparently consistent with a weakly correlated a Fermi liquid and \na Stoner enhancement of about a factor 2 to 3. \nThis argument again becomes inconclusive once the Fermi-liquid \nscenario is abandoned. In Fig.~\\ref{chi} we plot as function of $U\/W$ the \nuniform magnetic susceptibility in the normal phase calculated by DMFT. \nAfter an initial Stoner enhancement at small $U\/W$, the susceptibility flattens out \nand remains almost constant before a rapid growth which takes places extremely \nclose to the Mott transition. In the plateau region the susceptibility \nenhancement of a factor between 2 and 3 with respect to $U=0$, surprisingly \nclose to the experimental enhancement, covering the whole superconducting domain. \nPhysically, the origin of this susceptibility plateau for increasing $U$\nis quite instructive. It corresponds to the gradual crossover of the maximum \nspin available at each site from S=3\/2 in the Fermi liquid at small $U < U_{cp}$\n(Kondo screened AIM) to S=1\/2 near Mott and large $U > U_{cp}$, where Fermi liquid \nbehavior is lost. In essence, at large $U$ each molecule is effectively in a (dynamically) \nJT distorted state, where two out of three electrons are spin paired~\\cite{Manini94}. \nWe are thus led to conclude that the relatively weak observed enhancement of susceptibility \ndoes not correspond at all to a Stoner enhanced, weakly correlated Fermi liquid -- \nin fact quasiparticles do not even exist in most of the plateau region.\n \n\\begin{figure}\n\\includegraphics[width=7cm]{fig10.eps}\n\\caption{\\label{chi} Normal-phase uniform magnetic susceptibility $\\chi$ normalized to the \nnon-interacting value $\\chi_0$. For $U=0$, $\\chi\\simeq \\chi_0$, the difference being \nextremely small since $J\\ll W$. The plateau between $U\/W$ = 0.5 and 1 signals the \neffective crossover from a Fermi liquid with $S$= 3\/2 per site, to a non Fermi liquid\nwith $S$=1\/2 per site. }\n\\end{figure} \n\n\nFinally, we wish to address signatures of the strongly correlated scenario which \nwe expect should show up in important spectroscopies including \nthe tunneling I-V characteristics and angle resolved photoemission spectroscopy (ARPES). \nThis is initially embarrassing on two accounts. First, ARPES spectroscopies are $\\mathbf{k}$-vector resolved, whereas\nin DMFT we do not have access to any spatial structure. Second, tunneling spectroscopies \nare extremely well resolved near zero voltage, whereas our Lanczos method yields\na much poorer spectral function resolution in this region. \n\nLet us addressing tunneling first. Although Fig.~\\ref{DOS} refers to the impurity \nC$_{60}^{3-}$ molecule, we believe that similar features would remain after full DMFT \nself-consistency in the normal phase -- if we only had a better low-frequency numerical resolution \nthan we presently have. Hence, we suggest that tunneling I-V spectra of expanded \nfullerides be measured and examined, in order to bring out the expected rich structure \nof the kind sketched in Fig.~\\ref{DOS}. \n\nNext, let us consider photoelectron spectroscopy. Again according to the single-impurity \nanalysis~\\cite{DeleoPRL}, the imaginary part of the single-particle self-energy should \nbe finite and of order $T_+$ almost everywhere in the non-Fermi liquid normal phase above $T_c$. \nThis has the following implications for ARPES:\n\n\\begin{itemize}\n\n\\item the fulleride photoemission spectrum should show $t_{1u}$ bands dispersing \nin the Brillouin zone with nonzero bandwidth, governed by the energy scale $T_+$. \nThe value of $T_+$ decreases with increasing $U\/W$ (increasing expansion), from \n$W$ at $U$ = 0 to (larger than) $J$ at the Mott transition; \n\n\\item there should be a spectral peak broadening of the same order \nof magnitude $T_+$ as the apparent $\\mathbf{k}$-resolved band dispersion. In particular the broadening \nshould remain constant approaching the Fermi surface -- unlike a Fermi liquid phase \nwhere quasi-particle peaks become narrower and narrower.\n\n\\end{itemize}\n\\begin{figure}\n\\includegraphics[width=7cm]{fig11.eps}\n\\caption{\\label{Arpes} Simulated photoemission spectra elaborated form the \nDMFT result in the normal phase. The left panel refers to $U\/W = 1.1$ (corresponding \nto unexpanded fullerides), while the right panel is for $U\/W = 1.3$ (corresponding to \nexpanded fullerides).\n}\n\\end{figure} \n\nThe momentum-independence of the DMFT self-energy implies that, in this approximation, the \n$\\mathbf{k}$-modulation of the electronic dispersion is assumed to remain unaffected by interactions. \nWithin this assumption and without requiring too high frequency-accuracy, we can compute a toy $\\mathbf{k}$-resolved spectral function according to \n\n\\begin{equation}\n\\label{eq:arpes}\nA(k,\\omega)= -\\frac{1}{\\pi} \n\\mbox{Im} \\frac{1}{\\omega-\\varepsilon_{\\bf k} -\\Sigma_{DMFT}(\\omega)},\n\\end{equation}\n\nwhere $\\varepsilon_{\\bf k}$ is the non interacting dispersion and $\\Sigma_{DMFT}(\\omega)$ is the DMFT self-energy calculated with a finite \nnumber of baths. The effect \nof the local self-energy will be to change the effective bandwidth and to give \nrise to finite lifetime effects, even if the $\\mathbf{k}$-modulation of the dispersion \nis unrenormalized. For our Bethe lattice, there is no straightforward definition \nof momentum, and we remedy that by computing $A(\\varepsilon,\\omega)$, \nwhich corresponds to Eq. (\\ref{eq:arpes}) with $\\varepsilon_{\\bf k} \\to \\varepsilon$.\n\nIn Fig.~\\ref{Arpes} we show theoretical ARPES results for some choices of \n$\\varepsilon$ obtained by using the DMFT self-energy for temperatures above \n$T_c$, both for a value of $U\/W =1.1$ which lies close to the maximum of the \nsuperconducting dome, but still on the less correlated side, corresponding \nto unexpanded (or moderately expanded) fullerides, and for a value which \nlies in the downward branch of the dome ($U\/W = 1.3$), corresponding to a\nvery expanded fulleride. We note in both cases the existence of an incoherent \nlow-energy feature dispersing with a reduced but nonzero electron bandwidth \nof $0.1~W$. In the expanded case the pseudogap feature is clearly present.\n\nRecent photoemission spectra of K$_3$C$_{60}$~\\cite{Goldoni} indicate an\noverall dispersion bandwidth of about 160~meV, about a quarter of the bare \ncalculated bandwidth in the local-density approximation. The experimental spectral \npeak does not show the usual Fermi-liquid-like narrowing on approaching the Fermi \nlevel, a fact which is in agreement with our expectation for a non Fermi liquid \n(although a non expanded fulleride like K$_3$C$_{60}$ probably lies at the \nbeginning of the SCS dome, where deviations from Fermi liquid are not massive). \nExperimentally the peak does not appear to cross the Fermi level, and the \nintensity instead drops, suggestive of a pseudogap. Unfortunately the spectrum\nshows very strong vibronic effects, reflecting the retarded strong electron Jahn Teller \ncoupling. This aspect is not covered by our unretarded approximation, but it\nheavily affects the line shape and hampers the extraction of purely electronic features. \nTreatment of the vibronic effects, and a quantitative description\nof dispersion, will require abandoning in the future our approximation of infinitely \nfast phonon dynamics, as well as a possible extension to cluster extensions of DMFT which \nallow for different renormalizations of different momenta~\\cite{cdmft}.\n\n\n\\section{Conclusions}\nSummarizing, we addressed the apparently contradictory properties\nof expanded trivalent fullerides superconductors and insulators -- and to some extent \nof the whole family of fullerides -- and presented \na theoretical scenario emphasizing the role of strong electron correlations. \nThat is especially designed and appropriate for the more expanded\nmembers of the family, such as (NH$_3$)$_x$NaK$_2$C$_{60}$, Li$_3$C$_{60}$, \nCs$_{3-x}$K$_x$C$_{60}$ and Cs$_{3-x}$Rb$_x$C$_{60}$, and the recently discovered A15 $Cs_{3}C_{60}$ \nthat are near or past the Mott transition.\n\nOur model explains the dome-shaped increase and subsequent decrease of $T_c$ \nupon expansion of the lattice spacing in fullerides; the coexistence \nof metallic behavior and of Mott insulator features such as the large NMR spin \ngap in all fullerides, and the $S=1\/2$ spin in the insulating state (identified \nas a Mott-Jahn-Teller insulator). It explains why the $s$-wave $T_c$ can be \nas high as 40~K even though the Coulomb interaction strength is prohibitive, and\nwhy $T_c$ does not automatically decrease upon increase of $U\/W$. It also accounts \nfor more standard observations, such as regular specific heat jumps and moderately \nhigh spin susceptibilities, facts that were so far construed as evidence for conventional \nBCS superconductivity. \n\nBesides those listed in the previous Section, one can anticipate a number of additional experiments \nthat could provide ``smoking gun'' evidence for strongly correlated superconductivity\nin fullerides. The tunneling $I$-$V$ characteristics observable, e.g., by a scanning tunneling spectroscopy tip \nshould, in an expanded fulleride, develop the low energy\nfeatures typical of the Kondo impurity spectral function. The isotope effect\nupon carbon substitution should also behave very unconventionally, and eventually get \nsmaller as the superconducting dome is passed and the Mott transition is approached \nupon expansion. In this regime, as the quasiparticle bandwidth $ZW$ gradually falls below\nthe typical energy $\\hbar\\omega$ of an increasing fraction of the eight $H_g$ \nJahn Teller modes, the associated retardation effect should in fact disappear.\nThe expanded fullerides and related materials, clearly not enough\ninvestigated so far, deserve in our view the strongest experimental attention.\nThey combine elements that make them members at large of the high temperature \nsuperconductor family. They combine neighborhood of the Mott transition and predominance \nof strong electron correlations, with conventional elements such as electron-phonon $s$-wave \npairing, that are typical of BCS systems. Our study identifies a pseudogap and other \nfeatures in the IV tunneling spectrum, an increase of zero-frequency optical weight in the optical \nresponse of the superconducting phase, and the emergence of two separate energy scales \nin ARPES as the most urgent experimental undertakings that could confirm \nof falsify our claims.\n\n\n\\section*{Acknowledgments}\nE.T. thanks Kosmas Prassides, Andrea Goldoni, and Mauro Ricco' for information \nand discussion about fullerides and again K.P. for providing us with Fig. 1. \nM.C. acknowledges discussions with A. Toschi. Work in SISSA was sponsored by PRIN \nCofin 2006022847, as well as by INFM\/CNR ``Iniziativa trasversale calcolo parallelo''. \nWork in Rome was sponsored by PRIN Cofin 200522492. \n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nThe study of the kinetics of diffusion-controlled or \nballisticaly-controlled\nannihilation provides a nice playground on which the tools of \nnonequilibrium statistical mechanics can be tested ~\\cite{redner,born}. \nParticularly in low dimensions, the kinetics of these systems is \ngoverned by the fluctuations and any mean-field like (or Boltzmann-like) \napproximation will not give a satisfactory description of the properties \nof the system.\n\nA lot of efforts have be expended in studying the diffusion-controlled \ncase, both for the situation in which the reactants are uniformely \ndistributed in space~\\cite{wil} and\nfor the one in which they are initially spatially separated, leading to \nreaction-diffusion fronts~\\cite{racz,cordro}. In both cases, the long \ntime behavior of physical observables is characterized by power law. The \nassociated exponents have some universal properties~\\cite{born}.\n\nMore recently, the case of ballisticaly-controlled annihilation has been \ninvestigated for spatially homogeneous conditions~\\cite{reddue}. In \nparticular, exact results have been obtained for one-dimensional, single \nspecies models by Piasecki~\\cite{pia_uno} and Droz, Rey, Frachebourg \nand Piasecki~\\cite{pia_due,pia_tre}.\n\nIn these models, one considers point particles which move freely, with a\ngiven velocity. When two particles collide they instantaneously \nannihilate \neach other and disappear from the system. The system with only two \npossible\nvelocities $+c$ or $-c$ has been studied in a pioneering work by Elskens \nand\nFrisch~\\cite{ef}. The case of an arbitrary discrete velocity \ndistributions has been investigated by Droz, Rey, Frachebourg and \nPiasecki~\\cite{pia_tre}. An exact equation for the survival probability \nuntil the time $t$ of a particle moving with velocity $v$ was derived, \nallowing to compute the density and the time dependent velocity \ndistribution in the asymptotic regime $t\\to\\infty$.\nFor a symmetric three velocities distribution,\ndifferent kinetic regimes were found as a function of the initial \nfraction of particles at rest. \nSuch processes can model several physical situations as a recombination \nreaction in the gas \nphase or the fluorescence of laser excited gas atoms\nwith quenching on contact (the one-dimensional aspect can be obtained \nby working in a suitable porous media~\\cite{kopel}) or, \nthe annihilation of kink-antikink pairs in solid state \nphysics\\cite{buti}.\n\nA related but different class of problems is the case in which the two \nspecies (called $A$ and $B$), are initially separated in space and \nannihilate on contact. Such process can model the situation in which \nchemical species incorporated in a gel move ballistically under the \naction of a drift~\\cite{zrini}. As the two species cannot penetrate one \ninto the other (as they annihilate on contact), a well defined reaction \nfront is formed.\nWe aim at computing exactly the time dependent properties of this front.\n\nThe paper is organized as follows. In section two, the model is defined.\nIn section 3, we compute the average number of collisions that have \noccured before time $t$. An exact analytic expression for the \nprobability density to find the front at a given point and a given time \nis derived.\nIn section 4, it is shown that the front makes a random walk around its \n(time dependent) average value.\nPossible extensions of this work are discussed in the conclusion.\n\n\n\\section{The model}\n\nWe consider a one-dimensional system formed of two species of particles.\nInitially, particles $A$ are spatially randomly distributed in the \nregion $(-\\infty, 0)$\nand the $B$ ones are spatially randomly distributed in the region $(0, \n\\infty)$. The positions of the particles obey a Poissonian distribution. \nThe\ninitial linear average densities for the $A$ and $B$ particles are \nrespectively $\\mathrel{\\rho^{\\mathrm{A}}}$ and $\\mathrel{\\rho^{\\mathrm{B}}}$.\nThe velocities of each particle is an independent random variable \ntaking the value $\\pm c$ with even probability. Particles of the same \nkind suffer elastic collisions. When two particles $A$ and $B$ meet, \nthey annihilate.\nThus, practically the $A$ particles with velocity $-c$ and the $B$ \nparticles with velocity $+c$ move freely. Accordingly, the relevant part \nof the dynamics concerns the $A$ particles with velocity $+c$ and the \n$B$ ones with velocity $-c$.\n\nLet $(y_1,y_2,\\ldots,y_k,\\ldots)$ be the initial positions of the \n$A$ particles with velocity $+c$ and \n$(x_1,x_2,\\ldots,x_k,\\ldots)$ the initial positions of the $B$'s with \nvelocity $-c$ (see figure~\\ref{fig:one}). The relative velocity between \nthe $A$ and \n$B$ particles being $2c$, the pair of particles initially at $(x_k, \ny_k)$ will collide at time $t_k=(x_k-y_k)\/2$.\n\\begin{figure}[tbh]\n\\setlength{\\unitlength}{1cm}\n\\centerline{\n\\begin{picture}(12,3.5)(-6,-1)\n\\put(-5,0){\\line(1,0){10}}\n\\put(0,1){\\vector(-1,0){3}}\n\\put(-2.5,1.1){$A$ particles}\n\\put(0,1){\\vector(1,0){3}}\n\\put(2.5,1.1){\\makebox(0,0)[br]{$B$ particles}}\n\\put(0,0.9){\\line(0,1){0.2}}\n\\put(0,-0.2){\\line(0,1){0.4}}\n\\put(0,-0.3){\\makebox(0,0)[t]{$0$}}\n\\put(-0.5,0){\\circle*{0.1}}\n\\put(-0.5,-0.3){\\makebox(0,0)[t]{$y_1$}}\n\\put(-1.3,0){\\circle*{0.1}}\n\\put(-1.3,0.3){\\makebox(0,0)[b]{$y_2$}}\n\\put(-1.7,0){\\circle*{0.1}}\n\\put(-1.7,-0.3){\\makebox(0,0)[t]{$y_3$}}\n\\put(-2.7,-0.3){\\makebox(0,0)[t]{$\\cdots$}}\n\\put(-3.7,0){\\circle*{0.1}}\n\\put(-3.7,-0.3){\\makebox(0,0)[t]{$y_{k-1}$}}\n\\put(-4.4,0){\\circle*{0.1}}\n\\put(-4.4,0.3){\\makebox(0,0)[b]{$y_k$}}\n\\put(-5.1,0){\\makebox(0,0)[r]{$\\ldots$}}\n\\put(0.8,0){\\circle*{0.1}}\n\\put(0.8,-0.3){\\makebox(0,0)[t]{$x_1$}}\n\\put(1.1,0){\\circle*{0.1}}\n\\put(1.1,0.3){\\makebox(0,0)[b]{$x_2$}}\n\\put(1.8,0){\\circle*{0.1}}\n\\put(1.8,-0.3){\\makebox(0,0)[t]{$x_3$}}\n\\put(2.6,-0.3){\\makebox(0,0)[t]{$\\cdots$}}\n\\put(3.4,0){\\circle*{0.1}}\n\\put(3.4,-0.3){\\makebox(0,0)[t]{$x_{k-1}$}}\n\\put(4.2,0){\\circle*{0.1}}\n\\put(4.2,0.3){\\makebox(0,0)[b]{$x_k$}}\n\\put(5.1,0){\\makebox(0,0)[l]{$\\ldots$}}\n\\end{picture}}\n\\caption{Typical initial condition. The particles are labeled such that \n$y_k<\\cdotsy_{k+1}$. Hence averaging over the initial \nconditions, we find:\n\\begin{eqnarray}\nP(k;t) &=& \\langle\\theta\\left(2ct-[x_k-y_k]\\right)\n \\theta\\left([x_{k+1}-y_{k+1}]-2ct\\right)\\rangle \n \\nonumber\\\\\n &=& \\langle\\theta\\left(2ct-[x_k-y_k]\\right)\\rangle -\n \\langle\\theta\\left(2ct-[x_{k+1}-y_{k+1}]\\right)\\rangle\\, .\n \\label{eq:kcoll}\n\\end{eqnarray}\nwhere $\\theta$ is the usual Heaviside function and the brackets denote \nthe average over the initial conditions (see Appendix~\\ref{app:B}).\nThus:\n\\begin{equation}\nN(t)=\\sum_{k=1}^\\infty k P(k;t)\\, ,\n\\end{equation}\nwhich, using (\\ref{eq:kcoll}), reduces to \\begin{equation}\nN(t)=\\sum_{k=1}^\\infty \\langle\\theta\\left(2ct-[x_k-y_k]\\right)\\rangle\n\\, .\n\\end{equation}\nThus, using the results of Appendix~\\ref{app:B}, $N(t)$ reads:\n\\begin{eqnarray*}\nN(t)=\\sum_{k=1}^\\infty \\int_{-\\infty}^0 \\d y_k\\mathrel{\\rho^{\\mathrm{A}}} \\int_0^\\infty \\d x_k\n&\\mathrel{\\rho^{\\mathrm{B}}}& \n \\e^{-\\mathrel{\\rho^{\\mathrm{A}}}|y_k|}\\frac{\\left(|y_k|\\mathrel{\\rho^{\\mathrm{A}}}\\right)^{k-1}}{(k-1)!} \\\\\n&\\times&\n \\e^{-\\mathrel{\\rho^{\\mathrm{B}}} x_k}\\frac{\\left(x_k\\mathrel{\\rho^{\\mathrm{B}}}\\right)^{k-1}}{(k-1)!}\n \\theta\\left(2ct-[x_k-y_k]\\right)\\, .\n\\end{eqnarray*}\nNoticing that the sum over $k$ can be exchanged with the two integrals, \nwe \nfind (where $a=\\half(x_k-y_k)$ and $b=\\half(x_k+y_k)$)\n\\[\nN(t)=2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}} \\int_0^{ct} \\d a \\int_{-a}^a \\d b\n \\e^{-a(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})} \\e^{b(\\mathrel{\\rho^{\\mathrm{A}}}-\\mathrel{\\rho^{\\mathrm{B}}})}\n I_0\\left(2\\sqrt{\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}(a^2-b^2)}\\right)\\, ,\n\\]\n$I_0$ is the modified Bessel function. The integral over $b$ is \ncalculated \nin appendix~\\ref{app:A} and the integration on $a$ is straightforward. \nWe obtain\n\\begin{equation}\nN(t)=\\frac{2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}}{\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}}}\\left[ct - \\frac{1}{2(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})}\n \\left(1-\\e^{-2ct(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})}\\right)\\right]\\, .\\label{eq:nbcol}\n\\end{equation}\nIn terms of the dimensionless time\n\\[\n\\tau\\equiv ct(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})\n\\]\nand the parameter\n\\[\n\\lambda\\equiv\\frac{4\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}}{\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}}}\\, ,\n\\]\nwe have:\n\\begin{equation}\nN(\\tau)=\\frac{\\lambda}{2}\\left[\\tau -\n \\frac{1}{2}\\left(1-\\e^{-2\\tau}\\right)\\right]\\, .\n\\end{equation}\nAfter a transient regime, $N(\\tau)$ grows linearly with time.\n\nThe front position is defined as the point where the next collision \nwill \ntake place. Thus the probability density $\\mu(X;t)$ to find the front at \nthe point $X$, \nat time $t$ is given by\n\\begin{eqnarray}\n\\lefteqn{\\mu(X;t) =\n \\left\\langle\\delta\\left(X-\\frac{x_1+y_1}{2}\\right)\n \\theta(x_1-y_1-2ct)\\right\\rangle}\\nonumber \\\\\n&+& \\left\\langle\\sum_{k=1}^\\infty \\delta\\left(X \n-\n \\frac{x_{k+1}+y_{k+1}}{2}\\right)\n \\theta\\left(2ct-[x_k-y_k]\\right)\n \\theta\\left([x_{k+1}+y_{k+1}]-2ct\\right)\\right\\rangle \\nonumber \\\\\n&=& \n\\left\\langle\\delta\\left(X-\\frac{x_1+y_1}{2}\\right)\\right\\rangle\n \\label{eq:defmu} \\\\\n&+& \\left\\langle\\sum_{k=1}^\\infty \\left[\n \\delta\\left(X-\\frac{x_{k+1}+y_{k+1}}{2}\\right) -\n \\delta\\left(X-\\frac{x_k+y_k}{2}\\right) \\right]\n \\theta\\left(2ct-[x_k-y_k]\\right)\\right\\rangle\\, . \\nonumber\n\\end{eqnarray}\nThe first term in the right hand side reads simply\n\\begin{eqnarray}\n\\left\\langle\\delta\\left(X - \\frac{x_1+y_1}{2}\\right)\\right\\rangle\n&=& \\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}} \\int_0^\\infty \\d x_1 \\int_{-\\infty}^0 \\d y_1 \n \\e^{-\\mathrel{\\rho^{\\mathrm{B}}} x_1}\\e^{-\\mathrel{\\rho^{\\mathrm{A}}} y_1}\n \\delta\\left(X-\\frac{x_1+y_1}{2}\\right) \\nonumber\\\\\n&=& \\frac{2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}}{\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}}}\\left[\\theta(X)\\e^{-2\\mathrel{\\rho^{\\mathrm{B}}} X} +\n \\theta(-X)\\e^{2\\mathrel{\\rho^{\\mathrm{A}}} X}\\right]\\, .\n\\end{eqnarray}\n\nTo calculate the second term in the right hand side, we remark that one \nhas to\n average over the position $x_k$, $x_{k+1}$, $y_k$ and $y_{k+1}$. We \nshall then first average over particles $k+1$, to obtain an expression \ncontaining only the position $x_k$ and $y_k$. To do this, we need to \ncalculate the following integral\n\\[\nS \\equiv \\mathrel{\\rho^{\\mathrm{B}}}\\int_{x_k}^\\infty \\d x_{k+1} \\mathrel{\\rho^{\\mathrm{A}}}\\int_{-\\infty}^{y_k} \\d \ny_{k+1}\n\\delta\\left(X-\\frac{x_{k+1}+y_{k+1}}{2}\\right)\n\\e^{-\\mathrel{\\rho^{\\mathrm{B}}} x_{k+1}+\\mathrel{\\rho^{\\mathrm{A}}} y_{k+1}}\\, ,\n\\]\nputting $a=\\half(x_{k+1}-y_{k+1})$ and \n$b=\\half(x_{k+1}+y_{k+1})$, we have:\n\\begin{eqnarray*}\nS = 2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}\\int_{-\\infty}^\\infty \\d a \\int_{-\\infty}^\\infty \\d b \n\\theta(a+b-x_k)&\\theta&(a-b+y_k)\\delta(x-b) \\\\\n&\\times&\\e^{-(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})a+(\\mathrel{\\rho^{\\mathrm{A}}}-\\mathrel{\\rho^{\\mathrm{B}}})b}\n\\end{eqnarray*}\nThe integral over $b$ is straightforward, and the integral over $a$ \ngives after some algebra\n\\begin{eqnarray*}\nS = \\frac{2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}}{\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}}}\n&& \\left[\\rule{0mm}{5mm}\\theta\\left(\\frac{x_k+y_k}{2}-X\\right)\n \\e^{-(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})x_k+2\\mathrel{\\rho^{\\mathrm{A}}} X}\\right .\\\\\n&& +\\left.\\rule{0mm}{5mm}\\theta\\left(X-\\frac{x_k+y_k}{2}\\right)\n \\e^{(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})y_k-2\\mathrel{\\rho^{\\mathrm{B}}} X}\\right]\\, .\n\\end{eqnarray*}\nHence, the second term in the right hand side of eq.~(\\ref{eq:defmu}) \nbecomes\n\\begin{eqnarray*}\n\\sum_{k=1}^\\infty\\left\\langle\\left\\{\\rule{0mm}{6mm}\n \\frac{2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}}{\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}}}\\right.\\right.\n&& \\left[\\rule{0mm}{5mm}\\theta\\left(\\frac{x_k+y_k}{2}-X\\right)\n \\e^{-(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})x_k+2\\mathrel{\\rho^{\\mathrm{A}}} X}\\right .\\\\\n&& \n+\\left.\\rule{0mm}{5mm}\\theta\\left(X-\\frac{x_k+y_k}{2}\\right)\n \\e^{(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})y_k-2\\mathrel{\\rho^{\\mathrm{B}}} \nX}\\right]\\\\\n&& -\\left.\\left.\\rule{0mm}{6mm}\n \\delta\\left(X-\\frac{x_k+y_k}{2}\\right) \\right\\}\n \\theta\\left(2ct-[x_k-y_k]\\right)\\right\\rangle\\, .\n\\end{eqnarray*}\nNow using the relation (see appendix~\\ref{app:B})\n\\begin{eqnarray}\n\\lefteqn{\\sum_{k=1}^\\infty\\langle f(x_k+y_k,x_k-y_k) \\rangle}\n \\label{eq:sumav} \\\\\n&=& 2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}\\int_0^\\infty \\d a \\int_{-a}^a \\d b\n e^{-(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})a+(\\mathrel{\\rho^{\\mathrm{A}}}-\\mathrel{\\rho^{\\mathrm{B}}})b} f(2b,2a)\n I_0\\left(2\\sqrt{\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}(a^2-b^2)}\\right)\\, ,\\nonumber\n\\end{eqnarray}\none finally finds:\n\\begin{eqnarray}\n\\lefteqn{\\mu(X;t) = \\frac{2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}}{\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}}}\\left[\\theta(X)\n \\e^{-2\\mathrel{\\rho^{\\mathrm{B}}} X}+\\theta(-X)\\e^{2\\mathrel{\\rho^{\\mathrm{A}}} X}\\right]} \\\\\n &+& 2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}} \\int_0^{ct} \\d a \\int_{-a}^a \\d b\n \\e^{-(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})a+(\\mathrel{\\rho^{\\mathrm{A}}}-\\mathrel{\\rho^{\\mathrm{B}}})b}\n I_0\\left(2\\sqrt{\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}(a^2-b^2)}\\right)\n \\nonumber \\\\\n&\\times& \\left\\{-\\delta(X-b)+\\frac{2\\mathrel{\\rho^{\\mathrm{A}}}\\mathrel{\\rho^{\\mathrm{B}}}}{\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}}}\n \\left[\\theta(b-X)e^{-2\\mathrel{\\rho^{\\mathrm{A}}}(b-X)}+\n \\theta(X-b)e^{-2\\mathrel{\\rho^{\\mathrm{B}}}(X-b)}\\right]\\right\\}\\, . \\nonumber\n\\end{eqnarray}\nIn terms of the dimensionless variables $\\tau$ and $\\lambda$, and \n the dimensionless coordinate $\\xi\\equiv(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})X$, we can \nwrite\n\\begin{eqnarray}\n\\mu(\\xi;\\tau)\n&=& \\frac{\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}}}{2}\\e^{\\kappa\\xi} \\left\\{\\rule{0mm}{6mm}\n \\lambda\\e^{-|\\xi|}\\right. \\label{eq:mugen} \\\\\n&+& \\left.\\rule{0mm}{6mm}\\lambda\\int_0^\\tau \\d\\alpha \\e^{-\\alpha}\n \\int_{-\\alpha}^\\alpha \\d\\beta\n I_0\\left(\\sqrt{\\lambda(\\alpha^2-\\beta^2)}\\right)\n \\left[-\\delta(\\beta-\\xi)+\n \\frac{\\lambda}{2}\\e^{-|\\beta+\\xi|}\\right]\\right\\}\\, ,\\nonumber\n\\end{eqnarray}\nwith\n\\[\n\\kappa\\equiv\\frac{\\mathrel{\\rho^{\\mathrm{A}}}-\\mathrel{\\rho^{\\mathrm{B}}}}{\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}}}\\, .\n\\]\nNote that the function between \nthe curly brackets is an even function of $\\xi$; we shall call it \n$\\hat{\\mu}(\\xi;\\tau)$. The rest of the paper will be devoted to the \nstudy of this function.\n\n\n\\section{Results and interpretation}\\label{sec:results}\n\nStarting from~(\\ref{eq:mugen}) and using the integrals quoted in \nappendix~\\ref{app:A}, it is possible to verify that $\\mu(x;t)$ is \ncorrectly normalized to $1$. Moreover, we can also calculate its \ndifferent \nmoments. However for this purpose, it is simpler to consider the Fourier \ntransform of $\\mu$.\nLet us define the function ${\\tilde{\\mu}}(p;\\tau)$ as:\n\\begin{eqnarray}\n{\\tilde{\\mu}}(p;\\tau) &=& \\int_{-\\infty}^\\infty \\d X \n\\e^{-(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})pX}\n \\e^{-(\\mathrel{\\rho^{\\mathrm{A}}}-\\mathrel{\\rho^{\\mathrm{B}}})X} \\mu(X;\\tau) \\nonumber \\\\\n &=& \\frac{1}{2}\\int_{-\\infty}^\\infty \\d\\xi \\e^{-p\\xi}\n \\hat{\\mu}(\\xi;\\tau)\\, .\\label{eq:defmufou}\n\\end{eqnarray}\nHere again, this integral can be calculated using the formulas quoted in\nappendix~\\ref{app:A}. One finds:\n\\begin{equation}\n{\\tilde{\\mu}}(p;\\tau)=\\lambda\\e^{-\\tau}\\frac{1}{1-p^2}\\left[\n \\cosh\\left(\\sqrt{p^2+\\lambda}\\tau\\right)+\n \\frac{\\sinh\\left(\\sqrt{p^2+\\lambda}\\tau\\right)}\n {\\sqrt{p^2+\\lambda}}\\right]\\label{eq:mufou}\n\\end{equation}\nIt is now easy to calculate the different moment of $\\mu$, using the \nfact that\n\\begin{equation}\n\\langle\\xi^n\\rangle\n \\equiv \\int_{-\\infty}^\\infty \\d X\\left[(\\mathrel{\\rho^{\\mathrm{A}}}+\\mathrel{\\rho^{\\mathrm{B}}})X\\right]^n\\mu(X;t)\n =(-1)^n\\left[\\frac{\\partial^n}{\\partial p^n}{\\tilde{\\mu}}(p;\\tau)\n \\right]_{p=-\\kappa}\\, .\\label{eq:defmom}\n\\end{equation}\nNoticing that $\\kappa^2+\\lambda=1$, it is straightforward to check the \nnormalization. The first moment ${\\bar{\\xi}}(\\tau)=\\langle\\xi\\rangle$ is \ngiven by\n\\begin{equation}\n{\\bar{\\xi}}(\\tau) = \\kappa\\left[\\frac{2}{\\lambda}+\\tau-\\frac{1}{2}\n \\left(1-\\e^{-2\\tau}\\right)\\right]\n = \\frac{2\\kappa}{\\lambda}\\left[1+N(\\tau)\\right]\\, ,\n\\end{equation}\nand the variance is\n\\begin{eqnarray}\n\\lefteqn{\\xi_{\\mathrm{rms}}^2(\\tau) = \\langle\\xi^2\\rangle - \n\\langle\\xi\\rangle^2} \\\\\n&=& \\frac{4}{\\lambda^2} - \\frac{2}{\\lambda} + \\frac{3}{4} -\n \\frac{5\\lambda}{4} \n+\\lambda\\tau \n+\n \\left[\\lambda-\\frac{1}{2}-2\\tau(1-\\lambda)\\right]\\e^{-2\\tau} -\n \\frac{1-\\lambda}{4}\\e^{-4\\tau}\\, . \\nonumber\n\\end{eqnarray}\nIf $\\mathrel{\\rho^{\\mathrm{B}}} \\not= \\mathrel{\\rho^{\\mathrm{A}}}$, and in the long time limit, the front moves \nwith a velocity proportionnal to the density difference.\nThe case $\\lambda=1$ is of particular interest, because it corresponds \nto the symetric case $\\mathrel{\\rho^{\\mathrm{A}}}=\\mathrel{\\rho^{\\mathrm{B}}}$. As expected: \n\\begin{equation}\n{\\bar{\\xi}}(\\tau)\\equiv 0\n\\end{equation}\nbut\n\\begin{equation}\n\\xi^2_{\\mathrm{rms}}(\\tau) = 2\\left[1+N(\\tau)\\right]\\, .\n\\end{equation}\nIn the long time limit, we see that \n$\\xi_{\\mathrm{rms}}^2(\\tau)\\propto\\tau$:\nthe front move essentially as an unbiased random walker. Thus, we can \nexpect\n $\\mu$ to be essentially a gaussian. In fact, we shall prove that\nfor any $\\lambda$, the long time limit of $\\mu$ is a gaussian. \nStarting from the definition of an arbitrary moment of $\\mu$ \n(eq.~(\\ref{eq:defmom})), we may write for long time:\n\\begin{equation}\n\\langle\\xi^n\\rangle=(-1)^n\\lambda\\e^{-\\tau}\\left[\\frac{\\partial^n}\n{\\partial p^n}\\e^{\\tau\\sqrt{p^2+\\lambda}}\\right]_{p=-\\kappa}\n+ O(\\tau^{n-1})\\, ,\n\\end{equation}\nusing\n\\begin{equation}\n\\frac{\\partial^n}{\\partial p^n}\\e^{\\tau\\sqrt{p^2+\\lambda}}=\n\\left[\\frac{\\partial}{\\partial p}\\sqrt{p^2+\\lambda}\\right]^n\n\\left[\\frac{\\partial^n}{\\partial \\eta^n}\\e^{\\tau\\eta}\n\\right]_{\\eta=\\sqrt{p^2+\\lambda}} + O(\\tau^{n-1})\\, ,\n\\end{equation}\nwe obtain $\\langle\\xi^n\\rangle=(\\kappa\\tau)^n+O(\\tau^{n-1})$.\n\nIf $\\lambda<1$, $\\kappa\\neq 0$ and we know the leading term for every \nmoment of $\\mu$. If $\\lambda=1$, the situation is different and will be \nconsidered later.\n\nIt is known that for a gaussian distribution with a non vanishing mean \n${\\bar{\\xi}}(\\tau)$ \nand a variance $\\xi_{\\mathrm{rms}}^2(\\tau)$, one has:\n\\begin{equation}\n\\langle\\xi^n\\rangle_{\\mathrm{gauss}}=n!\\xi^n(\\tau)\n\\sum_{k=0}^{[\\frac{n}{2}]}\\frac{1}{(n-2k)!k!}\n\\left(\\frac{\\xi_{\\mathrm{rms}}^2(\\tau)}{2\\xi^2(\\tau)}\\right)^k\\, ,\n\\end{equation}\nIn the long time limit:\n\\begin{equation}\n\\frac{\\xi_{\\mathrm{rms}}^2(\\tau)}{2{\\bar{\\xi}}^2(\\tau)}\n=O\\left(\\frac{1}{\\tau}\\right)\\, ,\n\\end{equation}\nand $\\langle\\xi^n\\rangle_{\\mathrm{gauss}}=(\\kappa\\tau)^n+O(\\tau^{n-1})$.\n\nThus, when $\\lambda <1$ and for sufficiently long times (i.e. when the \ncorrections to the leading \nterm become negligible), the probabilty density $\\mu$ takes the gaussian \nform:\n\\begin{equation}\n\\mu(\\xi;\\tau)\\propto\n\\exp\\left[-\\frac{(\\xi-\\kappa\\tau)^2}{2\\lambda\\tau}\\right]\\, .\n\\end{equation}\nThe front wanders around its mean value with an amplitude of order of\nmagnitude $\\sqrt{\\lambda\\tau}$, like a biased random walker with a \ndiffusion coefficient of $\\lambda\/2$.\n\nThe $\\lambda=1$ case is different because the amplitude of the $\\tau^n$ \nterm vanishes. However, we can still prove that for large time, the \nprobability density $\\mu$ is a gaussian. More precisely, we shall give \nan upper and a lower bound, both of gaussian form. For this purpose, we \nshall come back to the Fourier transform ${\\tilde{\\mu}}(p;\\tau)$ defined \nby \neq.~(\\ref{eq:defmufou}). It can be shown (see appendix~\\ref{app:C}), \nthat the inverse Fourier transform of ${\\tilde{\\mu}}(p,\\tau)$ is given \nby\n\\begin{eqnarray}\n\\hat\\mu(\\xi;\\tau) &=&\\lambda\\e^{-|\\xi|-\\tau}\n \\left[\\cosh\\left(\\sqrt{1+\\lambda}\\tau\\right) +\n \\frac{\\sinh\\left(\\sqrt{1+\\lambda}\\tau\\right)}\n {\\sqrt{1+\\lambda}}\\right] \\label{eq:muborne}\\\\\n &-&\\lambda\\theta(\\tau-|\\xi|)\\e^{-\\tau}\n \\left\\{\\rule{0mm}{6mm}\n \\sinh\\left(\\sqrt{1+\\lambda}\\tau-|\\xi|\\right) +\n \\frac{\\cosh\\left(\\sqrt{1+\\lambda}\\tau-|\\xi|\\right)}\n {\\sqrt{1+\\lambda}}\\right.\\nonumber \\\\\n &-&\\left.\\rule{0mm}{6mm}\n 2\\int_0^{\\sqrt{\\lambda}}\\frac{\\d r}{\\pi}\n \\frac{\\cos(r\\xi)}{1+r^2}\\left[\n \\sinh\\left(\\sqrt{\\lambda-r^2}\\tau\\right)+\n \\frac{\\cosh\\left(\\sqrt{\\lambda-r^2}\\tau\\right)}\n {\\sqrt{\\lambda-r^2}}\\right]\\right\\}\\, .\\nonumber\n\\end{eqnarray}\nIt is now easy to find an upper and a lower bound for $\\hat\\mu$. Indeed, \n for $|\\xi|\\leq\\frac{\\pi}{2\\sqrt{\\lambda}}$, one has:\n\\[\n\\frac{\\cos(r\\xi)}{1+\\lambda}\\leq\\frac{\\cos(r\\xi)}{1+r^2}\\leq\\cos(r\\xi)\n\\]\nas $r\\in[0,\\sqrt{\\lambda}]$.\nThe remaining integral can be calculated exactly:\n\\begin{eqnarray}\n\\lefteqn{\\frac{2}{\\pi}\\int_0^{\\sqrt{\\lambda}} \\d r \n\\cos(r\\xi)\\left[\n \\sinh\\left(\\sqrt{\\lambda-r^2}\\tau\\right)+\n \\frac{\\cosh\\left(\\sqrt{\\lambda-r^2}\\tau\\right)}\n {\\sqrt{\\lambda-r^2}}\\right]}\\nonumber \\\\\n\\quad\n&=& \\lambda\\tau\\frac{I_1\\left(\\sqrt{\\lambda(\\tau^2-\\xi^2)}\\right)}\n {\\sqrt{\\lambda(\\tau^2-\\xi^2)}}+\n I_0\\left(\\sqrt{\\lambda(\\tau^2-\\xi^2)}\\right)\\, , \\label{eq:borne}\n\\end{eqnarray}\nyielding, then two bounds for $\\mu$. Two conclusions can then be\ndrawn. First for $\\lambda=1$, the asymptotic time behavior \nof~(\\ref{eq:borne}) is\n\\[\n\\e^\\tau\\frac{2}{\\sqrt{2\\pi\\tau}}\\e^{-\\xi^2\/2\\tau}\\, ,\n\\]\nand then, for $\\xi\\leq\\frac{\\pi}{2}$ and $\\tau$ large, we have:\n\\begin{equation}\n\\frac{1}{\\sqrt{2\\pi\\tau}}\\e^{-\\xi^2\/2\\tau}\\leq\\hat\\mu(\\xi;\\tau)\n\\leq\\frac{2}{\\sqrt{2\\pi\\tau}}\\e^{-\\xi^2\/2\\tau}\\, .\n\\end{equation}\nIn the long time limit, the density probability $\\mu$ becomes the \ndensity \nprobability of a random walk, with a diffusion coefficient of $\\half$.\n\nThe second consequence concerns the case $\\xi=0$, for any value of \n$\\lambda$: we have, in the long time limit:\n\\begin{eqnarray}\n\\hat\\mu(0;\\tau) &\\propto& \n\\frac{\\e^{(1-\\sqrt{\\lambda})\\tau}}{\\sqrt{\\tau}}\n \\quad\\mbox{if }\\lambda<1\\, , \\\\\n\\hat\\mu(0;\\tau) &\\propto& \\frac{1}{\\sqrt{\\tau}}\n \\quad\\mbox{if }\\lambda=1\\, .\n\\end{eqnarray}\nThese results are not really surprising. Indeed, it can be verified that \n$\\xi=0$ is \na maximum of $\\mu$, for $\\lambda=1$ but not for $\\lambda<1$.\n\nIt is useful to get an image of what is going on by plotting the \nfunction\n$\\mu$ for various value of $\\lambda$. However for this purpose, both \nequations~(\\ref{eq:mugen}) and~(\\ref{eq:muborne}) are inadequate for \ntechnical reasons. Equation~(\\ref{eq:mugen}) is not well suited\n because it contains a double integral which requires a lot of \nCPU time and equation ~(\\ref{eq:muborne}) because the integral diverge \nfor $r=\\sqrt\\lambda$, \nleading to numerical difficulties.\nTo avoid these problems, we have obtained a different \nintegral representation, starting from~(\\ref{eq:mugen}), and using the \nfollowing relation:\n\\begin{equation}\n\\left[\\frac{\\partial^2}{\\partial\\alpha^2} -\n \\frac{\\partial^2}{\\partial\\beta^2}\\right]\nI_0\\left(\\sqrt{\\lambda(\\tau^2-\\xi^2)}\\right) =\n\\lambda I_0\\left(\\sqrt{\\lambda(\\tau^2-\\xi^2)}\\right)\\label{eq:Ider} \\, .\n\\end{equation}\nThis relation follows directly from the integral representation for \n$I_0$\n\\[\nI_0\\left(\\sqrt{\\lambda(\\tau^2-\\xi^2)}\\right) =\n\\int_{\\mathrm{C}} \\frac{\\d z}{2\\mathrm{i}\\pi}\\frac{1}{z}\n\\exp\\left\\{\\frac{\\sqrt{\\lambda}}{2}\\left[(\\alpha+\\beta)z+\n \\frac{\\alpha-\\beta}{z}\\right]\\right\\}\\, ,\n\\]\nwhere $\\mathrm{C}$ is, for example, the circle of radius $1$ centered\nin $z=0$. Inserting~(\\ref{eq:Ider}) in~(\\ref{eq:mugen}), integrating \ntwice by part and using that\n\\[\n\\left[\\frac{\\partial^2}{\\partial\\alpha^2} - \n \\frac{\\partial^2}{\\partial\\beta^2}\\right]\n\\e^{-\\alpha-|\\beta-\\xi|} = 2\\e^{-\\alpha}\\delta(\\beta-\\xi)\\, ,\n\\]\nwe get\n\\begin{eqnarray*}\n\\hat\\mu(\\xi;\\tau) = \\lambda\\left\\{\\e^{-|\\xi|}\\rule{0mm}{10mm}\\right.\n + \\frac{1}{2}\\int_0^\\tau\n&\\d&\\alpha\n \\left[\\e^{-\\tau}\\left(\\e^{-|\\alpha-\\xi|}+\\e^{-|\\alpha+\\xi|}\\right)\n \\rule{0mm}{9mm}\\right. \\\\\n&\\times&\n \\left(\\lambda\\tau\n \\frac{I_1\\left(\\sqrt{\\lambda(\\tau^2-\\alpha^2)}\\right)}\n {\\sqrt{\\lambda(\\tau^2-\\alpha^2)}}+\n I_0\\left(\\sqrt{\\lambda(\\tau^2-\\alpha^2)}\\right)\\right)\\\\\n&-& \\left.\\rule{0mm}{10mm}\\left.\\rule{0mm}{9mm}\n 2\\theta(\\alpha-\\xi)e^{\\xi-2\\alpha} \n-\n 2\\theta(\\alpha+\\xi)e^{2\\alpha-\\xi}\\right]\\right\\}\\, .\n\\end{eqnarray*}\nPerforming the integral over the two theta functions, rearranging \nthe terms in the remaining integral and using the integrals of \nappendix~\\ref{app:A}, leads to:\n\\begin{eqnarray}\n\\hat\\mu(\\xi;\\tau)\n&=& \\lambda\\e^{-|\\xi|-\\tau}\n \\left[\\cosh\\left(\\sqrt{1+\\lambda}\\tau\\right) +\n \\frac{\\sinh\\left(\\sqrt{1+\\lambda}\\tau\\right)}\n {\\sqrt{1+\\lambda}}\\right] \\label{eq:muplot}\\\\\n&-& \\lambda\\theta(\\tau-|\\xi|)\\e^{-\\tau}\\left\\{\\rule{0mm}{10mm}\n \\sinh(\\tau-\\xi)\\right.\\nonumber \\\\\n&-& \\left.\\rule{0mm}{10mm}\n \\int_{|\\xi|}^{\\tau} \\d\\alpha\\sinh(\\alpha-|\\xi|)\\left[\n \\frac{I_1\\left(\\sqrt{\\lambda(\\tau^2-\\alpha^2)}\\right)}\n {\\sqrt{\\lambda(\\tau^2-\\alpha^2)}}+\n I_0\\left(\\sqrt{\\lambda(\\tau^2-\\alpha^2)}\\right)\\right]\\right\\}\\, .\n \\nonumber\n\\end{eqnarray}\nWe have not been able to simplify further this expression and in \nparticular to perform analytically the integration. However, this form \ncan be easily integrated numerically. For examples, the function \n$\\mu(\\xi;\\tau)$ is plotted on figure~\\ref{fig:two} and~\\ref{fig:three} \nfor respectively\n$\\lambda=1$ and for $\\lambda=0.5$, $\\mathrel{\\rho^{\\mathrm{A}}}>\\mathrel{\\rho^{\\mathrm{B}}}$.\n\\begin{figure}\n\\epsfysize=8cm\n\\centerline{\\epsfbox{mu_10.eps}}\n\\caption{Plot of $\\mu(\\xi;\\tau)$, for $\\mathrel{\\rho^{\\mathrm{A}}}=1$ and $\\mathrel{\\rho^{\\mathrm{B}}}=1$\n(i.e. $\\lambda=1$). \n\\label{fig:two}}\n\\end{figure}\n\\begin{figure}\n\\epsfysize=8cm\n\\centerline{\\epsfbox{mu_05.eps}}\n\\caption{Plot of $\\mu(\\xi;\\tau)$, for $\\mathrel{\\rho^{\\mathrm{A}}}=1+\\half\\protect\\sqrt{3}$ and \n$\\mathrel{\\rho^{\\mathrm{B}}}=1-\\half\\protect\\sqrt{3}$ (i.e. $\\lambda=0.5$). \n\\label{fig:three}}\n\\end{figure}\n\n\nOne sees that for $\\lambda=1$, the probability density $\\mu$ spreads out \nand approches a gaussian for long times. For the assymetric case \n$\\lambda=0.5$, one sees clearly the drift of the front which beecomes \nlinear in the long time regime.\n\n\\section{Conclusions}\\label{sec:concl}\n\nUsing a combinatorial analysis approach, we have been able to derive exact\nresults for the time dependent number of collisions, the probability\ndensity $\\mu(X;t)$ to find the front at point $X$ at time $t$ and\nits first two moments.\nAlthough simple, this nonequilibrium model leads to nontrivial behavior \nas far as the dynamics of the annihilation front is concerned. \n\nSeveral extensions of the present work are possible. One is to consider\na similar model in \nwhich the velocities are not restricted to be bimodal, but obey a more \ngeneral (discrete or continuous) distribution. Another interesting \naspect is given by the generalization of this model to higher \ndimensions. In this case, some $A$ particles can invade the right part \nof the system whithout being annihilated (and vice-versa for the $B$'s). \nA new definition of the front in term of the annihilation rate should be \nintroduced and different time dependent behavior may occur. These more \ncomplicated problems are under investigation.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\n\nElectrokinetic phenomena, which couple fluid flow, ion transport and\nelectric fields in electrolytes, are exploited in a variety of\nmicrofluidic technologies~\\cite{squires2005}. A solid surface in\ncontact with an electrolyte typically acquires a surface charge and\nforms an electric double layer composed of wall charge and a screening\nlayer of excess counterions. Electro-osmosis results from the action\nof an applied electric field on this double-layer charge, which sets\nthe bulk fluid into motion. Although the double layer is typically\nvery thin (nm) compared to the channel dimensions ($\\mu$m), it is the\nsource of long-range flow patterns. \n\n\nThe classical continuum theory of\nelectro-osmosis~\\cite{lyklema_book_vol2,hunter_book} suffices to\nunderstand most direct-current (DC) electro-osmotic flows, which are\nlinear in the applied voltage. DC electroosmosis results from\napplication of an external electric field parallel to the surface,\nwhich acts on equilibrium double-layer charge and induces motion in\nthe bulk, neutral liquid. The linear response is due to the assumption\nthat the applied field does not perturb the pre-existing surface\ncharge. When DC electroosmosis is used in microfluidic environments,\nelectrodes are placed at the ends of relatively long (cm) channels,\nand therefore large voltages (kV) are typically needed to generate\nsufficient fields.\n\nElectrokinetic phenomena can be exploited at much lower voltages and\nwith much greater precision by placing electrodes close together {\\em\n inside} the fluid channels and using alternating current (AC) to\ninhibit Faradaic reactions. Large electric fields (100 V\/cm) needed to\ninduce adequate flow can thus be generated with only a few volts by\nmicro-electrodes. The most studied example is alternating-current\nelectro-osmosis (ACEO) \\cite{encyclopedia_ACEO}, discovered by Ramos\net al. ~\\cite{ramos1999,green2000b,gonzalez2000,green2002}. When a\nperiodic array of interdigitated electrodes is placed inside the\nchannel, an AC signal applied to the electrodes can generate a steady\nset of micro-vortices. Ajdari showed that if geometrical asymmetry is\nintroduced, then these vortices can be rectified to drive net pumping\nin one direction over the array~\\cite{ajdari2000}. The original\nimplementation of this principle was a planar array of flat electrode\npairs of unequal widths and\ngaps~\\cite{brown2001,ramos2003,studer2004}, as shown in Figure\n~\\ref{fig:schematic}. Recent designs for ``3D ACEO'' pumps have\nachieved much faster flows (mm\/sec) using non-planar \nelectrodes, which more efficiently rectify opposing slip velocities in\na ``fluid conveyor\nbelt''~\\cite{bazant2006,urbanski2006,urbanski2007}. Waves of\nvoltage can also pump fluids over electrode arrays by traveling-wave\nelectro-osmosis (TWEO)~\\cite{cahill2004,ramos2004}.\n\n\\begin{figure}\n\\centering\n \\epsfig{file=figure1.eps,width=2.8in,clip=} \n \\caption{\\label{fig:schematic} One period of an asymmetric planar\n array of microelectrodes in an AC electro-osmotic (ACEO)\n pump. This geometry has been studied experimentally by several\n groups~\\cite{brown2001,studer2004,urbanski2006,microTAS2007} and\n is the subject of the present theoretical study. The dimensions\n are; $W1=4.2~ \\mathrm{\\mu m}$ , $W2=25.7~\\mathrm{\\mu m}$, $G1=4.5~\n \\mathrm{\\mu m}$ , and $G2=15.6 ~\\mathrm{\\mu m}$, using the\n notation of Ref.~ \\cite{olesen2006}. }\n\\end{figure}\n\nThese effects exemplify the fundamental nonlinear electrokinetic\nphenomenon~\\cite{encyclopedia_nonlinear} of induced-charge\nelectro-osmosis (ICEO)~\\cite{iceo2004a,iceo2004b,levitan2005}, which\nalso has other applications in microfluidics. The key source of\nnonlinearity is that the diffuse double-layer charge is {\\em induced}\nnear a polarizable surface by the applied electric field, which then\nacts on it to drive nonlinear electro-osmotic flow. Although ICEO\nflows were first described in the 1980s in the context of colloid\nscience~\\cite{murtsovkin1996,gamayunov1986}, they are finding new\napplications in microfluidic devices. As in the case of ACEO pumps,\nbroken spatial symmetries in ICEO flows can also be used to pump\nfluids around polarizable microstructures and to manipulate polarizable\nparticles~\\cite{iceo2004a,squires2006,yariv2006,velev}.\n\nThe current theory of ICEO flow~\\cite{iceo2004b}, including\nACEO~\\cite{gonzalez2000,olesen2006} and\nTWEO~\\cite{cahill2004,ramos2004} microfluidic pumps, is based on the\nclassical electrokinetic\nequations~\\cite{lyklema_book_vol2,hunter_book}. This century-old model\ncomprises the Poisson-Nernst-Planck (PNP) equations of ion transport\ncoupled to the Navier-Stokes equations of viscous fluid flow via\nelectrostatic Maxwell stresses. The crucial modification is in the\nboundary conditions, relaxing the assumption of fixed surface charge\n(or constant zeta potential) to allow for significant electrostatic\npolarizability of the surface (e.g. fixed surface potential). The\nelectrokinetic equations themselves, however, have not been\nquestioned, until very recently~\\cite{large}, which provides the\nmotivation for this work.\n\nICEO flows typically involve {\\it large voltages} induced across the\ndouble layer, greatly exceeding the thermal equilibrium voltage,\n$kT\/e$ ($=25$ mV at room temperature). In particular, experimental\nACEO electro-array\npumps~\\cite{brown2001,studer2004,urbanski2006,urbanski2007,microTAS2007},\nthe subject of our study, usually apply several volts to the double\nlayer. Such large voltages of order $100 kT\/e$ inevitably lead to the\nbreakdown of the PNP equations by violating the fundamental assumption\nof an ideal dilute solution~\\cite{kilic2007a,kilic2007b}. The\nlimitations of dilute solution theory for double-layer structure are\nwell known, and many attempts have been made to incorporate effects\nsuch as steric exclusion and electrostatic correlations. (See\nRef.~\\cite{kilic2007a} for a recent review.) ICEO flow, however,\nraises new issues due to its more extreme, dynamical\ncontext~\\cite{large}. Even if the bulk solution is dilute, the surface\nbecomes so highly charged by the induced voltage that counterions\nbecome crowded in the double layer, while simultaneously driving\ntangential fluid flow.\n\n\\begin{figure}\n\\centering\n\\epsfig{file=figure2.eps,width=2.8in}\n\\caption{Pumping velocity vs. frequency for the planar ACEO pump in\n Figure~\\ref{fig:schematic} from experiments and simulations.\n {\\it Points:} Experimental data\n at 3 Vpp in 0.03 mM KCl from Ref.~\\cite{microTAS2007} (taken by\n J. P. Urbanski); {\\it Solid curve: } Linear Debye-H\\\"uckel-Stern\n model of the diffuse layer with a constant compact layer capacitance\n $C_s = \\varepsilon\/(\\delta \\lambda_D)$, where $\\delta=0.5$ is chosen\n to fit the experimental peak frequency. {\\it Dashed curve: }\n Nonlinear Gouy-Chapman-Stern (Poisson-Boltzmann) model of the\n diffuse layer also with $\\delta=0.5$. }\n \\label{fig:jpdata}\n\\end{figure}\n\nThe need for improved models is also indicated by some unexplained\nfeatures of ICEO experiments. In this paper, we focus on the tendency\nfor planar ACEO pumps to reverse at high\nvoltage~\\cite{studer2004,chang2004,wu2006,urbanski2006,microTAS2007},\nwhich naturally poses problems for microfluidic applications, if it\ncannot be reliably predicted. The standard model of ACEO based on the\nclassical electrokinetic equations predicts a single peak in the flow\nrate versus frequency at the characteristic ``RC'' charging frequency\nof the electrodes, in both the linear~\\cite{ajdari2000} and\nnonlinear~\\cite{olesen2006,olesen_thesis} regimes. Flow reversal has\nbeen observed at high voltage ($> 2$ V) and high frequency (10-100\nkHz) in ACEO pumping of dilute\nKCl~\\cite{studer2004,chang2004,microTAS2007} and deionized\nwater~\\cite{urbanski2006} with $10 \\mu$m scale electrode\narrays. \n\nRepresentative examples of the observed and predicted frequency\nresponse are shown in Figure \\ref{fig:jpdata} for the pump geometry of\nFigure~\\ref{fig:schematic}. Here we compare experimental data taken at\n3.0 V $= 120 k T\/e$ in an 0.03 mM KCl solution over gold electrodes\n~\\cite{microTAS2007} to simulations with two commonly used models: one\nassumes the linear Debye-Huckel (DH) model of the double\nlayer~\\cite{ramos1999,gonzalez2000,ajdari2000}, while the other adopts\na non-linear Poisson-Boltzmann (PB) model based on the Gouy-Chapman\nsolution~\\cite{olesen2006}. Following prior\nwork~\\cite{ajdari2000,green2002,levitan2005}, a constant Stern-layer\ncapacitance is added in order to better {\\em fit} the experimental data by\nmatching the peak frequency in the DH model. \nBoth models can reasonably fit the data in terms of the \npeak in the forward pumping, though \nthis becomes more difficult for PB at higher voltages \\cite{olesen2006}. \nEach model predicts a single peak of forward flow,\nwhile the experiments show the pump reversing above 10 kHz. The\ninability of these classical models to capture, even qualitatively,\nthis high frequency flow reversal motivates us to consider modified\ndouble-layer models.\n\nThe reversal of ACEO pumps has been considered in several previous\nstudies. Flow reversal was first attributed to Faradaic\nreactions~\\cite{chang2004}, and this view of ``Faradaic charging''\npersists in recent work on DC-biased AC\nelectro-osmosis~\\cite{wu2007}. Simulations of ACEO pumps with\nButler-Volmer reaction kinetics, however, have failed to predict the\nobserved flow, especially at high\nfrequency~\\cite{olesen2006,olesen_thesis}. In existing models, a weak\nflow reversal due to reactions can only be observed at low frequency\n(far below the maximum) and for certain sets of\nparameters~\\cite{ajdari2000,olesen2006,olesen_thesis}. This is\nconsistent with the observation of gas bubbles from electrolysis at\nlow frequency and high voltage in dilute KCl with gold\nelectrodes~\\cite{studer2004}. Recently, higher resolution\nmeasurements of the same pump design with platinum electrodes have\nrevealed very weak ($<10 \\mu$m\/sec) reverse ACEO flow at low frequency\n($< 20$ kHz) and low voltage ($<1.5$ V) and have demonstrated the\nimportance of Faradaic currents through {\\it in situ} impedance\nspectroscopy~\\cite{gregersen2007}. Although the theory predicts\nsimilar flow reversal, it is not in quantitative agreement and does\nnot predict the concentration dependence.\n \nIn addition to Faradaic reactions, various other nonlinear effects in\ndilute solution theory also dominate at low frequencies: The\ndifferential capacitance of the diffuse layer~\\cite{lyklema_book_vol2}\ndiverges, which causes the RC charging time to grow exponentially with\nvoltage~\\cite{olesen2006}, and salt adsorption and tangential\nconduction by the diffuse layer are coupled to (much slower) bulk\ndiffusion~\\cite{bazant2004,chu2006,chu2007}. All of these effects have\nrecently been incorporated in simulations of ACEO at large voltages\nusing the classical electrokinetic equations, but high-frequency flow\nreversal was not observed~\\cite{olesen2006,olesen_thesis}. Another\npossible source of flow reversal is AC electrothermal\nflow~\\cite{gonzalez2006}, which can lead to reverse pumping in\nexperiments with planar electrode arrays~\\cite{wu2007}, but under\ndifferent conditions of much higher salt concentration ($>$ 1 M),\nvoltage ($> 10$ V) and higher frequency ($> 100$ kHz) than ACEO ($<\n0.01$ M, $< 5$ V, $< 100$ kHz). To date, none of the above effects\nhave been able to explain the experimental data for ACEO pumps.\n\nIn this paper, we propose steric effects of ion crowding in the double\nlayer as a possible cause of the observed flow reversal at high\nfrequency. Using simple modifications of the electrokinetic equations\nto account for finite ion sizes, we are able to predict a frequency\nresponse of ACEO pumps that is very similar to that observed in\nexperiments. Although our results do not provide a complete\nquantitative theory, we will demonstrate that accounting for steric\neffects can have an important impact on the theory of ACEO.\n\n\n\\section{Steric effects in a thin double layer}\n\\label{sec:steric} \n\nIn the dilute-solution theory of\nelectrolytes~\\cite{lyklema_book_vol2,hunter_book}, the concentration of\neach ionic species in the diffuse part of the electric double layer at\na charged surface is in thermal equilibrium with a Boltzmann\ndistribution,\n\\begin{equation}\n c_i(x)=c_0 \\mathrm{exp}{\\left( \\frac{-z_ie\\phi(x)}{kT}\\right)}, \n\\label{eq:cminus}\n\\end{equation}\nwhere $k$ is Boltzmann's constant, $T$ the temperature, $e$ the\nelectron charge, $z_i$ the valence, $c_0$ the concentration in the\nneutral bulk electrolyte just outside the double layer, and $\\phi$ is\nthe (mean) electrostatic potential relative to the bulk ($x\\gg\n\\lambda_D$, where $\\lambda_D$ is the Debye-H\\\"uckel screening\nlength). In classical linear electrokinetic phenomena, the\ndiffuse-layer voltage drop $\\phi(0)$ is set by chemical equilibrium at\nthe surface and is thus typically comparable to the thermal voltage\n$kT\/e$. At an electrode driving ACEO flow, however, much larger\ndiffuse-layer voltages $\\phi(0) \\gg kT\/e$ are {\\it induced} by a\napplied voltage of order $100 kT\/e$. Under these conditions, it is\neasy to see that Boltzmann equilibrium (\\ref{eq:cminus}) breaks down\nby predicting diverging concentrations of counterions at the surface,\nof order $e^{100}c_0$. \n\nThe assumption of a dilute solution is thus incompatible with a large\napplied voltage. If we assume ions have a characteristic length scale,\n$a$, then the corresponding cutoff concentration $a^{-3}$ for the\nbreakdown of dilute-solution theory is reached at relatively low\nvoltage, even if the bulk salt solution is very dilute. For example,\nfor $c_0=10^{-5}$ M, $z=1$, and $a=3$ \\AA (including a solvation\nshell), the cutoff concentration is reached at $0.33$ V. To account\nfor the excess ions (at typical equilibrium voltages), Stern\npostulated a compact layer of solvated ions of finite size on the\nsurface~\\cite{bockris_book}, which carries most of the double-layer\nvoltage as the diffuse-layer capacitance diverges. Such an intrinsic\n``surface capacitance'' is also invoked in models of ICEO flows, where\nit can also include the effect of a thin dielectric coating on a metal\nsurface~\\cite{ajdari2000,bazant2004}. It seems unlikely, however,\nthat an atomically thin Stern or coating layer could withstand several\nVolts, e.g. since dielectric breakdown occurs in most materials\n(including water) in fields of order 10 MV\/m $= 0.01$\nV\/nm~\\cite{jones1995}. Instead, it seems the diffuse layer must carry\na substantial voltage $\\phi \\gg kT\/e$ in ACEO experiments, which\ncauses the ions to become highly crowded near the surface.\n\n\\begin{figure}\n(a)\\includegraphics[width=2.8in]{figure3a.eps}\n(b)\\includegraphics[width=2.8in]{figure3b.eps}\n\\caption{\\label{fig:C} (a) Schematic of the \n equilibrium distribution of counterions\n near a positively charged surface \n taking into account a minimum ion spacing $a$ for large applied\n voltages. As the voltage increases, the width of the condensed\n layer will increase causing a decrease the capacitance. \n (b) The voltage dependence\n of the capacitance of the diffuse \n layer for PB and MPB at different values of $\\nu$; \n where $\\nu = 2a^3 c_0$ is the\n bulk volume fraction of ions. }\n\\end{figure}\n\n\n\nA variety of ``modified Poisson-Boltzmann equations'' (MPB) have been\nproposed to describe equilibrium ion profiles near a charged wall, as\nreviewed in Ref.~\\cite{kilic2007a}. To capture ion crowding effects\nacross a wide range of voltages, we employ the simplest possible MPB\nmodel~\\cite{kilic2007a,kralj-iglic1996,borukhov1997,biesheuvel2007},\nfirst proposed by Bikerman~\\cite{bikerman1942},\nwhich is a continuum approximation of the entropy of ions on a lattice\nof size $a$. As shown in Fig.~\\ref{fig:C}(a), when a large voltage is\napplied, the counterion concentration exhibits a smooth transition\nfrom an outer PB profile to a condensed layer at $c=c_{max}=a^{-3}$\nnear the surface.\n\n\nIn ACEO, charging dynamics are very important, so the double-layer\ncapacitance is an important property. Dilute-solution theory predicts\nthat the differential capacitance, $C_D$, diverges with the voltage\nas,\n\\begin{equation}\n C_{D}(\\Psi_D)=\\frac{\\varepsilon_b}{\\lambda_{D}}\\cosh\\left(\\frac{ze\\Psi_{D}}{2kT},\n \\right) \\label{eq:cdpb}\n\\end{equation}\nwhere $\\Psi_D$ is the voltage applied across the double layer. For a\nconcentrated solution described by Bikerman's model, we predict the\nopposite trend~\\cite{kilic2007a},\n\\begin{equation}\nC_{D}^{\\nu }=\n\\frac{\\frac{\\varepsilon}{\\lambda_{D}}\\sinh(\\frac{ze\\Psi_{D}}{kT})}\n{[1+2\\nu\\sinh^2\\left(\\frac{ze\\Psi_{D}}{2kT}\\right)]\n\\sqrt{\\frac{2}{\\nu} \\mathrm{ln}[1+2\\nu\\sinh^2\\left(\\frac{ze\\Psi_{D}}{2kT}\\right)]}},\n\\label{eq:cdnu}%\n\\end{equation}\nwhere $\\nu=2 c_0 a^3$ is the bulk volume fraction of\nions~\\cite{kilic2007a}. (Note that the same formula was derived by\nKornyshev~\\cite{kornyshev2007} in the context of ionic liquids around\nthe same time that Kilic et al.~\\cite{kilic2007a} derived it in the\npresent context of electrolytes.) \n\nAs shown in Fig.~\\ref{fig:C}(b), with this model the capacitance\nreaches a maximum near the critical voltage then {\\it decreases} at\nlarger voltages because the capacitor width grows due to steric\neffects. The same result also follows from a simpler model using PB\ntheory outside a condensed layer of variable width at the maximum\ndensity ~\\cite{kilic2007a}. We will show that this general decrease\nin capacitance with voltage that can drastically impact the charging\ndynamics in ACEO flows.\n\nAs Biesheuvel {\\it et al.} point\nout~\\cite{biesheuvel2005,biesheuvel2007}, Bikerman's latticed-based\napproach can underestimate the true steric effect to a large degree,\nand more accurate models of the entropy of a hard-sphere liquid, such\nas the Carnahan-Starling (CS) equation of state or various extensions\nto multicomponent mixtures~\\cite{carnahan1969,lue1999}, can be used\ninstead. The CS equation is not tractable for a closed-form\nanalytical expression for double-layer capacitance, but we can easily\nobtain it numerically as a function of the voltage. For the details,\nsee the Appendix.\n\nAlthough various mean-field models, such as CS, are very accurate for\nequilibrium properties of the hard-sphere gas and liquid at low to\nmoderate volume fractions, $\\Phi < 0.55$, none can claim to accurately\ndescribe the packing limit at large volume fractions, which would\ngenerally lead to the decrease in double-layer capacitance at large\nvoltage~\\cite{kilic2007a}. For example, these models predict a\ndiverging chemical potential only at $\\Phi=1$, and not and not at the\nmathematical packing bound ($\\Phi=0.74$, corresponding to the\nface-centered cubic lattice) or the more relevant bound of the jamming\ntransition (around $\\Phi=0.63$ for hard spheres\n~\\cite{torquato00}). Even in simple model systems in equilibrium, the\njamming transition is poorly understood\n~\\cite{liu98,liu02,ohern02,ohern03,donev04}, so none of these models\ncan possibly be expected to accurately capture the dynamics of jamming\nof hydrated ions very close to a surface driven by a large AC voltage,\nwhile also producing ACEO flow. Of course, under such extreme\nconditions, the mean-field approximation for both steric and\nelectrostatic interactions should also be questioned.\n\nDue to all of these complexities, following Bazant {\\it et\n al.}~\\cite{large}, we will focus on the simplest approximation of\nBikerman's model in order to explore the qualitative impact of steric\neffects on ACEO flow. It will turn out that this very simple model\nsuffices to predict high frequency flow reversal of ACEO pumps. We\nwill then repeat some calculations with the CS model, which improves\nthe agreement with experiments somewhat, but does not alter the\nqualitative predictions of Bikerman's model.\n\n\n\\section{Model for ACEO pumps}\n\n\nIn our theoretical study of ACEO pumping, we consider the standard\nexperimental geometry of planar, asymmetric electrode-array pumps\nshown in Fig. \\ref{fig:schematic} due to Brown {\\it et\n al.}~\\cite{brown2001}. The specific dimensions are selected to\nmatch the experiments Studer \\textit{et. al} \\cite{studer2004} and\nUrbanski {\\it et al.} ~\\cite{urbanski2006,microTAS2007}. The basic\nquestion of ACEO flow in such pumps has been thoroughly addressed by\nOlesen {\\em et. al} ~\\cite{olesen2006,olesen_thesis} using the\nclassical electrokinetic equations of dilute solution\ntheory~\\cite{lyklema_book_vol2,hunter_book}. The significant\ndifference between this work and Ref. \\cite{olesen2006} is our\napplication of the new physical model for the double-layer\ncapacitance, taking into account steric effects at large\nvoltage. Unlike Ref.~\\cite{olesen_thesis}, we also ignore ``strongly\nnonlinear'' effects associated with salt uptake by a highly charged\ndiffuse layer and resulting bulk diffusion~\\cite{bazant2004,chu2006},\nsince these phenomena are greatly reduced by steric\neffects~\\cite{kilic2007a,kilic2007b}.\n\nIn the bulk fluid outside the\nelectric double layers, we assume electroneutrality with a constant\nelectrical conductivity, $\\sigma$. Under these assumptions, the\nelectric potential in the fluid bulk, $\\phi$, satisfies Laplace's\nequation, which is equivalent to Ohm's law for a constant bulk\nresistance. This corresponds to the ``weakly nonlinear'' regime where\nnonlinear circuit models still\nhold~\\cite{bazant2004,chu2006,olesen_thesis}. For boundary conditions\non the potential, we assume that there is no normal Ohmic current on\nthe sections of channel substrate at $y=0$ and the upper boundary at\n$y=H$.\n\nFor the boundary condition over the electrodes, we must incorporate a\nmodel of the thin electric double layers that will form. These double\nlayers are so thin (10-100 nm), that when we apply the boundary\ncondition on $\\phi$ at $y=0$, we assume this is outside the double\nlayer and not at the electrode surface. The local charge stored in\nthe double layer, $q(x)$, changes in time due to the Ohmic current.\nOur double layer model provides the voltage difference between the\nfluid bulk and the externally applied voltage at the electrode,\n$V_{ext}$, based on the amount of charge stored. Thus the boundary\ncondition for the electric potential in the bulk is dynamically\ndetermined. While Faradaic current certainly play a role in ACEO, we\nneglect it in our model to focus on steric effects alone, which have\nnot previously been considered. The role of reactions in classical\ndouble-layer models in ACEO was studied by Olesen {\\em\n et. al}~\\cite{olesen2006} and shown not to play a role in\nhigh-frequency flow reversal.\n\n\nAs in prior work~\\cite{ajdari2000,olesen2006}, we introduce the following\nscales for the nondimensionalization of the governing equations,\n\\[\n \\quad[x,y] = L,\\quad[t] = \\frac{L C_o}{\\sigma}, \\quad[\\phi ]=\\frac{k_B T}{z e},\\quad [q]=\\frac{k T }{z e }C_o.\n\\]\nNote that $[t]$ is the ``RC time'' for the equivalent\ncircuit~\\cite{ramos1999,bazant2004}, where the characteristic Debye\ncapacitance is $C_o = \\frac{\\varepsilon}{ \\lambda_D}$ where\n$\\varepsilon$ is the permittivity of the electrolyte and $\\lambda_D$\nis the Debye length. The characteristic length scale for computing\nthe bulk resistance, $L\/\\sigma$, is the electrode spacing $L=G_1$. We\nassume no Stern layer with our new double layer model, though its\ninclusion would be trivial and would not significantly alter our\nresults.\n\nThe dimensionless problem reduces to Laplace's equation for the\npotential in the fluid bulk,\n\\begin{equation}\n\\nabla^2 \\phi = 0. \n\\label{eq:lap}\n\\end{equation}\nThe boundary condition on the \nsubstrate (the entire upper wall at $y=H\/L$ and the insulating \nregions along $y=0$ between the electrodes) is\n\\begin{equation}\n\\frac{\\partial \\phi}{\\partial y}=0, \n\\label{eq:bc1}\n\\end{equation}\nThe boundary condition over the electrode double layer\nis given by, \n\\begin{equation}\n \\phi = V_{ext} - \\Psi_D(q), \n\\label{eq:bc2}\n\\end{equation}\nwhere $\\Psi_D(q)$ is the \ntotal potential drop across the double layer. The functional \nrelationship between the double layer voltage drop and charge\nis given by the \nmodel from Kilic \\textit{et al.} \\cite{kilic2007a},\n\\begin{equation}\n\\Psi_D = - 2~ \\mbox{sign}(q)~ \\mbox{sinh}^{-1} \\left( \\sqrt{ \\frac{1}{2 \\nu} \\left( \\mbox{e}^{\\frac{q^2 \\nu}{2}}-1 \\right)} \\right),\n\\label{eq:q_full}\n\\end{equation}\nwhile the local charge is dynamically determined by the ``RC''\ncondition,\n\\begin{equation}\n\\frac{dq}{dt} = \\left. \\frac{\\partial \\phi}{\\partial y} \\right|_{y=0}.\n\\label{eq:charge}\n\\end{equation}\n\n\n\nThe closed set of equations (\\ref{eq:lap})-(\\ref{eq:charge}) form our\nelectrical model. Note that in the limit of very small $\\nu$ in\nEq. \\ref{eq:q_full} we obtain the classic result from\nPoisson-Boltzmann, $\\Psi_D = -2~ \\mbox{sinh}\\left(\\frac{q}{2}\n\\right).$ In the limit of small potential and small $\\nu$ we obtain\n$\\Psi_D = -q,$ the linear Debye-Huckel model. The later two limits\nhave been extensively studied \\cite{olesen2006}; it is the impact of\nEq \\ref{eq:q_full} on the flow that is the focus of this paper.\n\nOnce the electrical problem is solved, we compute the electroosmotic\nslip by the Helmholtz-Smoluchowski relationship,\n\\begin{equation}\nU = \\Psi_D \\frac{\\partial \\phi(y=0)}{\\partial x}.\n\\label{eq:velocity}\n\\end{equation}\nand ignore any modifications at large voltage, e.g. due to crowding\neffects~\\cite{large}, in order to focus on steric effects in the\ndouble-layer charging dynamics. Velocity is made dimensionless by the\nnatural scale~\\cite{ramos1999,iceo2004a},\n\\[\n[U] = \\frac{\\varepsilon V^2}{\\eta L}\n\\]\nwhere $V$ is the magnitude of the applied voltage and $\\eta$ is the\nviscosity. An interesting fact, exploited by previous\nauthors~\\cite{ajdari2000,gonzalez2000,olesen2006}, is that we can\ncompute the net pumping without actually solving the flow field in the\nchannel, but by taking the time and spatial average of equation\n\\ref{eq:velocity} along the boundary $y=0$. Further details of this\nderivation can be found elsewhere, e.g. in Olesen {\\em\n et. al}~\\cite{olesen2006}.\n\nDue to the irregular electrode geometry and the non-linear charging\nprocesses, we must resort to solving Equations\n(\\ref{eq:lap})-(\\ref{eq:charge}) numerically. We couple the solution\nof Laplace's equation to the dynamic boundary condition for the\ncharge. The solution is integrated forward in time until an adequate\nsteady state is reached and the net flow is computed by averaging\nEq. \\ref{eq:velocity} in time and space. Details on the numerical\nmethod can be found in the supplementary material \\footnote{ See EPAPS\n Document No. XXXX for details of the numerical method. For more\n information on EPAPS, see http:\/\/www.aip.org\/pubservs\/epaps.html }.\nThe simulations were confirmed by an independent method using the \ncommercial software package COMSOL Multiphysics. \n\n\n\n\n\n\\section{Results}\n\n\\begin{figure}\n\\centering\n\\epsfig{file=figure4.eps,width=2.8in}\n\\caption{Average slip velocity as a function of frequency for three different\n models at 2.5 Volts. The solid blue curve is for $\\nu=0.01$ with\n the complete model, the black dash-dot line assumes the linear\n Debye-Huckel relationship $q=-\\Psi_D$, the red dashed curve is for\n $\\nu=0$ with a Stern layer capacitance of $\\delta=0.1$, Only the\n model with steric effects predicts flow in the negative\n direction. Points A, B, and C will be referred to in\n Fig. \\ref{fig:flow}. }\n \\label{fig:models}\n\\end{figure}\n\nWe start by comparing our results with the steric model for the double\nlayer, Eq. (\\ref{eq:q_full}), to simulations of the dilute-solution\nmodels used in all prior work. In Fig. \\ref{fig:models}, we compare\nthe flow velocity as a function of frequency for the classic linear\nDebye-Huckel model valid only when the applied voltage is much less\nthat $k T\/e$, the Poisson-Boltzmann non-linear capacitance which is\nrecovered from our model when $\\nu=0$, and the model which accounts\nfor steric effects, Eq. \\ref{eq:q_full}. At a voltage of $100\nkT\/e=2.5$ V, we see that the model which accounts for steric effect\nshows a high frequency flow reversal. As we note in\nRef.~\\cite{large}, to our knowledge, these results represent the first\ndemonstration of a mathematical model that can predict such reversal\nof flow in ACEO pumps.\n\n\\begin{figure*}\n\\centering\n \\epsfig{file=figure5.eps,width=5.8in,clip=} \n \\caption{\\label{fig:mechanism} Basic mechanism for flow in ACEO\n devices. Here we consider a symmetric array of electrodes under a\n suddenly applied voltage. In (a) we show electric field lines at\n the instance the field is turned on. The electric field points\n from the positive to the negative electrode. The location on the\n electrodes at $y=0$ is shown. In (b) we show the charge density\n over the electrodes at $1\/10$ the charging time $t_c$, at the\n charging time, and at $10$ times the charging time. In (c) we show\n streamlines of the electric field at the charging time. In (d) we\n show electric field lines at the charging time. The arrows denote\n the flow direction. }\n\\end{figure*}\n\nIn order to understand the physical mechanism of flow reversal, we\nmust first understand why there is flow at all. The basic mechanism\nfor flow in ACEO can be best understood when considering a\ngeometrically symmetric electrode array, as in Fig\n\\ref{fig:mechanism}, subjected to a suddenly applied DC voltage\n(without any Faradaic current) ~\\cite{encyclopedia_ACEO}. When the\nvoltage is initially switched on, the electric field lines point from\nthe positive to the negative electrodes as shown in Fig\n\\ref{fig:mechanism}(a). Since there is initially no double layer\ncharge, the field lines are perpendicular at the electrode surface. A\ncurrent will begin to flow through the electrolyte with ions migrating\nalong the electric field lines and the double layers over the\nelectrode will begin to charge as shown in Fig \\ref{fig:mechanism}(b).\nThe electric field is strongest at the electrode edges and the\nresulting high current quickly charges the double layer at that\nlocation. On the order of the characteristic charging time, $t_c$,\n(as computed using the gap length scale) the electrode edges will be\nsignificantly screened. In Fig \\ref{fig:mechanism}(b) we show the\ncharge at $y=0$ at three instances in time.\n\nAs the double layers charge up, the electric field lines in the fluid\nbulk are diverted around the screened electrode edges, resulting in a\ntangential component of the electric field. This is shown in Fig.\n\\ref{fig:mechanism}(c) where we can see the field lines are no longer\nperpendicular over the entire electrode, especially near the electrode\nedges. The tangential electric field exerts a force on the fluid in\nthe charged double layer, causing flow. Given the direction of the\nfield and sign of the charge, the flow always goes from the edge of\nthe electrodes, inward. The flow streamlines and flow direction are\nshown in Fig. \\ref{fig:mechanism} (d). If the applied voltage is\nmaintained, eventually the entire electrode becomes fully screened\nthere is no electric field in the bulk and the flow will stop. The\nabove argument about the direction of flow did not depend upon the\nsign of the applied voltage, reversing the sign causes the same flow\npattern. Therefore, if an AC voltage is applied at a frequency\ncorresponding to the charging time, time averaged vortex flow can be\nmaintained.\n\n\n\n\n\\begin{figure*}\n\\centering\n \\epsfig{file=figure6.eps,width=5.8in,clip=} \n \\caption{\\label{fig:flow} Streamlines (a,c,e) and slip velocity\n (b,d,f) for three different cases. The first row is for $V=100$,\n $\\nu=0.01$, and $\\omega = 0.7$. This frequency corresponds to the\n peak in the forward flow for the steric model. In (a) we show the\n streamlines. In (b) we show the time-averaged slip velocity (solid\n line) and the cumulative spatial integral of this velocity (the\n dashed curve). The second row, figures (c) and (d) is for $V=100$,\n $\\nu=0.01$, and $\\omega = 3$. This frequency corresponds to the\n peak in the reversed flow. The third row, figures (e) and (f) is\n for $V=100$ and $\\omega = 3$ with the linear Debye-Huckel model\n $\\Psi_D=-q$. }\n\\end{figure*}\n\nAjdari introduced the idea of rectifying this symmetric flow pattern\nto create an ACEO pump by introducing various forms of geometrical\nasymmetry within each spatial period of the electrode array\n\\cite{ajdari2000}. This principle was first implemented by Brown {\\it\n et al.} using flat, co-planar electrodes of different widths and\ngaps~\\cite{brown2001}, as shown in Fig. \\ref{fig:schematic}, and this\ngeneral design has been studied extensively by several\ngroups~\\cite{ramos2003,mpholo2003,studer2004,olesen2006,\nolesen_thesis,gregersen2007,urbanski2006,microTAS2007}.\nIn this geometry, the symmetry breaking causes the flow to go from\nleft to right by upsetting the delicate balance between leftward and\nrightward slip patterns. The fact that pumping comes from a relatively\nsmall bias of opposing, co-planar surface flows makes this design\nsusceptible to flow reversal if nonlinearities in charging dynamics\ntip the balance the other way, as we demonstrate below. Current\ndesigns using stepped, three-dimensional electrodes are much faster\nand more robust against flow reversal, since the opposing slip\npatterns work together to drive pumping in one direction, as a ``fluid\nconveyor\nbelt''~\\cite{bazant2006,urbanski2006,urbanski2007,burch_preprint}. The\nsensitivity to flow reversal in the planar pump, however, allows us a\nbetter opportunity to explore subtle nonlinearites in double-layer\ncharging dynamics, which is why we have chosen it for our theoretical\nstudy.\n\n\n\nIn Fig. \\ref{fig:flow} we demonstrate how the geometrical asymmetry\nresults in forward flow when driven at the charging time and how\nsteric effects can produce high frequency reversal. In\nFig. \\ref{fig:flow}(a) and (b) we show the time averaged streamlines\n(a) and time average slip velocity (b) for the ACEO pump driven at the\ncharging time, corresponding to point A in Fig. \\ref{fig:models}. We\nnotice a few features of the slip profiles. First, the magnitude of\nvelocity over the small electrode is greater than the velocity over\nthe large electrode. This difference is due to the fact that the {\\em\n total} charge contained on both electrodes must be the same, the\nlocal charge density (and thus electroosmotic slip) is always higher\non the small electrode. Second, we notice slip profiles over either\nelectrode is asymmetric with the greatest slip found at the edges at\nthe small gap. This effect is simply due to high electric field in the\nsmall gap which can exert a larger force. Finally we see the\nasymmetry in the slip is greater for the large electrode than the\nsmall electrode, an effect due to the geometry. The result for this\ngeometry and frequency is the net flow goes in the positive $x$\ndirection. The relative contributions to the flow is easily seen by\nthe cumulative spatial integral of the time averaged slip; the dotted\nline in Fig. \\ref{fig:flow} (b).\n\n\nIn Fig. \\ref{fig:flow} (c) and (d) we show the behavior at the\napproximate peak of the reverse flow in Fig. \\ref{fig:models},\ncorresponding to point B in Fig. \\ref{fig:models}. The same generic\nfeatures of the slip profiles are the same as in (b). However, the\nrelative contributions have changed at this frequency. The reason for\nthe reversal can be explained by the change in charging time for the\ndifferent sized electrodes. When steric effects are included, we must\nremember that the capacitance decreases as more charge is stored\ninside the double layer. Since the charge density over the small\nelectrode is much greater, the capacitance of the small electrode's\ndouble layer is {\\em decreased} at high voltage due to the steric\neffects. The lower capacitance means that the charging time of the\nsmall electrode is decreased relative to the large electrode. This\nchange in charging time means that at high frequency and high voltage\nthe small electrode has more time to sufficiently charge to generate\nflow.\n\nThis idea is illustrated in Fig. \\ref{fig:flow} (c) and (d). The\ncumulative spatial integral of the time-averaged slip shows that the\nnegative flow contribution from the small electrode is nearly twice\nthe positive contribution. At this high frequency the large electrode\nis not very highly charged and its forward flow component adds only a\nsmall amount to the net. That this reversal is caused by the change\nin the charging dynamics can be seen in Fig. \\ref{fig:flow} (e) and\n(f), which shows the same condition as (c) and (d) only with the\nlinear Debye-Huckel model; point C in Fig. \\ref{fig:models}. Here,\nsince the large electrode has the same charging time as the small\nelectrode, its forward component is significant to provide a net\npositive flow. It is clear from these pictures that with the\ninclusion of steric effects, the direction and magnitude of flow in\nACEO pumps is determined by a delicate balance influenced {\\em both}\nby the geometry and the charging dynamics.\n\n\\begin{figure*}\n\\centering\n\\epsfig{file=figure7.eps,width=5.8in,clip=}\n\\caption{Flow response as a function of voltage for various values of\n $\\nu$, the steric parameter. We show values of $\\nu=0.01$ (a),\n $\\nu=0.001$ (b), $\\nu=0.0001$ (c), $\\nu=0$ (d), The contours show\n constant levels of total flow rate. The heavy line shows the $U=0$\n contour and the region above this line denotes the region where flow\n reversal occurs. }\n \\label{fig:steric_nu}\n\\end{figure*}\n\nUp to now, we have taken the volume fraction, $\\nu$, to be a constant.\nTo better see the effect of $\\nu$ on the predicted flow, in\nFig. \\ref{fig:steric_nu} we show contours of the flow rate as a\nfunction of voltage and frequency at different values of $\\nu$. For a\nfixed physical value of $a$, different values of $\\nu$ represent how\nthe flow would change as a function of concentration (in dimensionless\nterms). The heavy lines denote the region where the flow is negative,\nindicating a reversal of flow. We see that the frequency response of\nthe flow does not change with voltage at low voltage, as would be\npredicted by the linear Debye-Huckel model. At low voltage, we see\nonly a single peak in the forward flow at the characteristic\nfrequency. As the voltage is increased the nonlinearity begins to set\nin which causes the peak of the forward flow to move toward lower\nfrequency as predicted by dilute solution theory. As the voltage is\nfurther increased the steric effects begin to become important with\nthe reversal eventually setting in.\n\n\n\n\nThe change in frequency of the maximum forward flow at higher voltage\nas $\\nu$ changes is easily understood. As $\\nu$ increases, the value\nof the maximum value of the capacitance decreases, shifting the peak\nresponse to higher frequencies. For a given $\\nu$, as voltage\ncontinues to increase, the continual decrease in the double layer\ncapacitance can be seen in the shift of the maximum of the positive\nflow to higher frequency at high voltage. When there are no steric\neffects (or Stern layer) the capacitance diverges, which can be seen\nin Figure \\ref{fig:steric_nu}(d). The inclusion of the constant Stern\nlayer capacitance causes this decrease to limit at the constant Stern\nlayer capacitance \\cite{olesen2006}.\n\nQualitatively, our model can predict a flow \nresponse that looks very similar to experimental\nresults, however the quantitative ability of our model is still problematic. \nIn Fig. \\ref{fig:studer}(a) we show the model predicting\nthe experimental data of\nStuder {\\it et al.} \\cite{studer2004}. In Fig. \\ref{fig:studer} we treat\nthe parameter $\\nu$ as a fitting parameter to approximate the\ndata of Figure 6 in Studer's experimental work \\cite{studer2004}. We\nfind a value of $\\nu=0.01$\nprovides a reasonable fit to their experimental data.\nHowever, given a concentration of 0.1 mM KCl, this value \ncorresponds to $a=4.4$ nm, about\nan order of magnitude larger than we would expect based on the hydrated\nradius of the ion. \nThe fact their experiments \\cite{studer2004}\nonly predicts reverse flow at high voltage (their Fig. 7) \nis consistent with the low\nfrequency regime of their work being unaccessible due to\nelectrochemical reactions and degradation of the electrode (see their Fig. 4).\n\n\\begin{figure}\n\\centering\n(a)\\includegraphics[width=2.8in]{figure8a.eps}\n(b)\\includegraphics[width=2.8in]{figure8b.eps}\n\n \\caption{Flow as a function of frequency at different voltages. The voltages are\n $V = 0.5,1,1.5,2.9,3.6,4.5,$ and $6$ Volts RMS. The voltages are\n selected to match Figure 6 in Studer {\\em et al.} \\cite{studer2004}.\n The data is presented in dimensional form in order to \n compare the frequency response. \n The flow velocity is corrected for the hydraulic resistance of the\n experimental flow loop. In (a), using the Bikerman model, we \n used $a=4.4$ nm. \n In (b), using the Carnahan-Starling model, we \n used $a=2.2$ nm. The general shape of both frequency response\n curves have many of the experimentally observed features. \n }\n \\label{fig:studer}\n\\end{figure}\n\n\n\nWe considered whether the Carnahan-Starling (CS) form of the excess\nchemical potential would provide a more reasonable value of $a$. In\norder to include the CS equation in the model we numerically solve for\nthe relationship of $\\Psi_D(q)$ to replace Eq. \\ref{eq:q_full} in our\nmodel. We follow the work of Bieshuvel and Soestbergen\n\\cite{biesheuvel2007} to obtain this relationship numerically, as\ndescribed in the Appendix. The result of using the CS equation in our\nACEO model is shown in Fig. \\ref{fig:studer} (b), and it is\nqualitatively similar to our prior results with the Bikerman\nmodel. The CS steric model has somewhat different quantitative\nbehavior but still requires an ion size of approximately $a=2.2$ nm to\nfit the experimental data. While the magnitude of the flow reversal\nchanges with the CS equation, the use of a more accurate model for the\nexcess chemical potential still require a relatively large ion size.\n\n\n\\section{Discussion}\nIn order to fit the experimental data, the ACEO models modified for\nsteric effects imply an ion length scale, $a$, almost an order of\nmagnitude larger than we would expect. A factor of 10 in ion size\nresults in a factor of 1000 in $\\nu$. This difference means that while\nthe response in Fig. \\ref{fig:steric_nu}(a) is closer to what is\nobserved, it is Fig. \\ref{fig:steric_nu} (c) that is predicted based\non a reasonable hydrated radius for the ions. The current model\ncannot yet predict the experimental data from first principles. It is\nimportant to note that the large value of $\\nu$ is needed not to\npredict flow reversal at all, but to predict flow reversal at the\nfrequencies observed in experiments.\n\nAs discussed in section ~\\ref{sec:steric}, the mean-field models of\nsteric effects which we consider are unlikely to accurately capture\nthe dynamical formation of a condensed layer on the surface. One\nreason is the difficulty in describing the jamming transition, as ions\nand solvent molecules are squashed onto the surface by large,\ntime-dependent normal electric fields. Under such extreme conditions,\nwe should not expect more than qualitative trends from such\nmodels. Another possibly crucial feature missing in these models is\nthe interaction with the surface, which is assumed to be a\nmathematically flat, homogeneously charged, perfectly rigid wall, as\nsketched in Fig. 2. Since the condensed layer of (presumably solvated)\nions is at the scale of only a few molecular lengths, there must be a\nstrong coupling with nanoscale roughness on the surface. This may have\nthe effect of increasing the effective ion size, as we have inferred,\nsince what matters is not the ion density in the solution, but the\ndensity very close to the surface, which is reduced by atomic\nasperities. The viscosity may also diverge near the surface as ions\nbecome jammed among surface inhomogeneities~\\cite{large}. It would be\ninteresting to do experiments with atomically flat electrodes\n(e.g. carbon graphene) or ones with controlled nanoroughness to test\nthis hypothesis.\n\n\nIt is also possible that changes in permittivity, $\\varepsilon$,\nwithin the double layer may also effect the frequency response.\nElectrochemists infer for the Stern layer that $\\varepsilon$ may be\n$1\/10$ of $\\varepsilon$ in the bulk ~\\cite{bockris_book}. If it were\ntrue that $\\varepsilon$ is reduced by $1\/10$ in the condensed layer,\nthe overall double layer capacitance would be decreased. Such change\nwould cause the peaks in the frequency response to shift to higher\nfrequency for the same value of $\\nu$. If we make the assumption that\n$\\varepsilon$ is reduced by a factor of ten in the {\\em entire} double\nlayer, then an ion size of approximately 1 nm is needed to predict\nStuder {\\em et al.}'s \\cite{studer2004} data with the CS equation. If\n$\\varepsilon$ is reduced only in the condensed layer, the true effect\non the frequency response would be less than predicted with the simple\nassumption. Therefore, it seems likely that accounting for changes in\n$\\varepsilon$ would reduce the needed ion size to correlate the data,\nhowever at most this would reduce the ion size to 1 nm, still much\nlarger than expected.\n\n\nA further complication is that the mean-field approximation breaks\ndown when ion spacings approach the Bjerrum length, $l_B =\n(ze)^2\/4\\pi\\varepsilon kT$, which is 7~ \\AA ~ for bulk water and\nmonovalent ions ($z=1$). Again, if it were true that $\\varepsilon$ is\nreduced by 1\/10 in the condensed double layer, then $l_B=7$ nm.\nCorrelation effects on electro-osmotic flow (which to our knowledge\nhave never been studied) could be very significant at large voltages,\neven in dilute bulk solutions~\\cite{large}. \n\n\n\nWhile much work remains, we have highlighted an interesting and\npractically important application where dilute solution theory clearly\nbreaks down. We have demonstrated that accounting for steric effects\nin double-layer models can have a dramatic impact on the predicted\nflow in ACEO pumps. Because of the large value of $a$ needed to fit\nthe experimental data, we cannot claim that our model definitively\nexplains high-frequency flow reversal in ACEO pumps, but it is the\nonly plausible explanation to date. The prediction is also quite\nrobust: we have found that flow reversal can be predicted by any model\nwhich accounts steric effects by the generic feature that the double\nlayer capacitance decreases with applied voltage. \n\nPhysically, it seems that this generic effect of decreasing\ncapacitance must be true for any model of the double layer considering\nfinite sized ions~\\cite{kilic2007a}. As the voltage increases, more\nions are packed in the double layer. Once the concentration of ions\nbecomes so high that the packing density is reached, the ions would\nhave to pile up forming a condensed layer whose thickness would\nincrease with voltage. As the double layer thickness grows, the\ncapacitance must decrease. We believe this general effect is true\nregardless of the details of the model, and we have shown that it has\nmajor implications for ACEO pumping and, by extension, all nonlinear\nICEO flows. Steric effects are clearly among the new effects which\nmust be considered in developing more accurate theories of ICEO flow\nat large applied voltages~\\cite{large}.\n\n\\section*{Acknowledgments}\nThis work was supported in part by the National Science Foundation\nunder contract DMS-0707641 and by the U.S. Army through the Institute for\nSoldier Nanotechnologies, under Contract DAAD-19-02-0002 with the\nU.S. Army Research Office. MZB also acknowledges support from ESPCI\nthrough the Joliot Chair.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzdlqz b/data_all_eng_slimpj/shuffled/split2/finalzzdlqz new file mode 100644 index 0000000000000000000000000000000000000000..408feec13171f9dc1fb51e964ba8528fedf89362 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzdlqz @@ -0,0 +1,5 @@ +{"text":"\\section*{Introduction}\n\t\t\n\t\tHurwitz numbers are a classical topic in combinatorics and algebraic\n\t\tgeometry; they date back to the work by Adolf Hurwtiz \\cite{AdHurwitz}\n\t\t(1890). A classical definition of the Hurwitz number $h_{m,\\lambda}$\n\t\twhere $m$ is a non-negative integer and $\\lambda=(\\lambda_1 \\DT\\ge\n\t\t\\lambda_s)$ is a partition, is: $h_{m,\\lambda}$ is the number of\n\t\tsequences of $m$ transpositions $(\\sigma_1 \\DT, \\sigma_m)$ such that\n\t\ttheir product $\\sigma_1 \\DT\\circ \\sigma_m$ belongs to a conjugacy\n\t\tclass $C_\\lambda \\in S_n$ in the symmetric group $S_n$. By a classical\n\t\tresult \\cite{GouldJacks}, generating function of the Hurwitz numbers\n\t\tsatisfies a second order parabolic PDE called the cut-and-join\n\t\tequation. Its right-hand side has Schur polynomials as eigenvectors,\n\t\tso the Hurwitz numbers can be expressed via them\n\t\t\\cite{LandoKazarian}. The result by A.\\,Okounkov about Toda lattice\n\t\t\\cite{Okounkov} implies that the same generating function is a\n\t\tsolution of the KP hierarchy (the proof uses tools coming from\n\t\tmathematical physics, namely the boson-fermion\n\t\tcorrespondence). Finally, in a recent work \\cite{TwistedHurwitz} by\n\t\tYu.Burman and the author Hurwitz numbers are shown to solve a problem\n\t\tin low-dimensional topology: they are the numbers of ways to glue $m$\n\t\tribbons to a collection of $n$ disks so as to obtain a surface with a\n\t\tprescribed structure of the boundary.\n\t\t\n\t\tIn this paper we show that there is a parallel theory of Hurwitz\n\t\tnumbers for the reflection groups of series B and D. Section\n\t\t\\ref{Sec:BD} is the starting point where we recall the classical\n\t\tembedding of $B_n$ and $D_n$ into the symmetric group $S_{2n}$; the\n\t\treflections are mapped to transpositions and products of pairs of\n\t\ttranspositions commuting with the involution $\\t = (1,n+1) (2,n+2)\n\t\t\\dots (n,2n)$ (for the group $D_n$, only pairs of transpositions are\n\t\tused). Also we recall a description of the conjugacy classes in these\n\t\tgroups, taken from \\cite{Carter} (see Proposition \\ref{Pp:ConjClass};\n\t\tthe classes are generally indexed by pairs of partitions $\\lambda$ and\n\t\t$\\mu$ such that $\\lmod \\lambda\\rmod + \\lmod \\mu\\rmod = n$).\n\t\t\n\t\tSection \\ref{Sec:HurwNum} gives the actual definition of Hurwitz\n\t\tnumbers for the groups $B_n$ and $D_n$: it is similar to the classical\n\t\tone with reflections instead of transpositions. For the B series, we\n\t\tcount reflections of two classes (transpositions and pairs of\n\t\ttranspositions) separately.\n\t\t\n\t\tIn Sections \\ref{Sec:Expl} and \\ref{Sec:CJD} we study how the\n\t\tmultiplication by a reflection affects conjugacy classes; this gives\n\t\tus an expression for the cut-and-join operator for the groups of\n\t\tseries B and D. After a suitable change of variables we see that the\n\t\toperator obtained for the B group is actually a tensor square of the\n\t\tclassical cut-and-join. The cut-and-join for the D group is a direct\n\t\tsum of the cut-and-join for B and the classical cut-and-join\n\t\trescaled. This leads to an expression of the Hurwitz numbers for B and\n\t\tD via Schur polynomials.\n\t\t\n\t\tIn Section \\ref{SSec:Ribbon} we provides a model for B-Hurwitz\n\t\tnumbers involving ribbon decomposition of surfaces equipped with an\n\t\tinvolution. Finally, in Section \\ref{SSec:Fermion} we apply the\n\t\tboson-fermion correspondence to the tensor square of the cut-and-join\n\t\tdescribed above; it allows one to show that the generating function of\n\t\tthe B-Hurwitz numbers is a $2$-parameter family of $\\tau$-functions of\n\t\tthe KP hierarchy, independently in $2$ variables, and to prove a\n\t\tsimilar result for D.\n\t\t\n\t\t\\subsection*{Acknowledgements}\n\t\t\n\t\tThe author wants to thank his advisor Yurii Burman for many many\n\t\tfruitful discussions and constant attention to this work.\n\t\t\n\t\tThe research was partially funded by the HSE University Basic Research\n\t\tProgram and by the International Laboratory of Cluster Geometry NRU\n\t\tHSE (RF Government grant, ag.~No.~075-15-2021-608 dated 08.06.2021)\n\t\t\n\t\t\\section{Reflection groups $B_n$ and $D_n$}\\label{Sec:BD}\n\t\t\\subsection{Definition and embedding into $S_{2n}$}\n\t\t\n\t\tLet $B_n$ be a rank $n$ finite reflection group with the mirrors\n\t\t$x_i-x_j = 0$, $x_i+x_j = 0$ ($1 \\le i < j \\le n$) and $x_i = 0$ ($1\n\t\t\\le i \\le n$); the corresponding reflections are denoted $s_{ij}^+$,\n\t\t$s_{ij}^-$, and $\\ell_i$, respectively (see \\cite{Humphreys} for details\n\t\tand standard notation). The subgroup $D_n \\subset B_n$ is generated by\n\t\t$s_{ij}^{\\pm}$ only.\n\t\t\n\t\tConsider an involution $\\t \\in S_{2n}$ without fixed points. All such\n\t\tinvolutions are conjugate, and the results that follow do not depend\n\t\ton a particulate choice; still, fix $\\t \\stackrel{\\mbox{\\tiny def}}{=} (1,n+1)(2,n+2) \\dots\n\t\t(n,2n)$ (that is, $\\t(i) = i+n \\bmod 2n$ for all $i$).\n\t\t\n\t\tGroups $B_n$ and $D_n$ admit embedding into the permutation group\n\t\t$S_{2n}$. Namely,\n\t\t\n\t\t\\begin{proposition}\\label{Pp:Embed}\n\t\t\tThere exists an embedding $\\Phi: B_n \\to S_{2n}$ such that\n\t\t\t$\\Phi(s_{ij}^+) = (ij)(i+n,j+n)$, $\\Phi(s_{ij}^-) = (i,j+n)(i+n,j)$\n\t\t\tand $\\Phi(\\ell_i) = (i,i+n)$. The image $\\Phi(B_n)$ is the normalizer\n\t\t\tof $\\tau$:\n\t\t\t\\begin{equation*}\n\t\t\t\t\\Phi(B_n) = \\{\\sigma \\in S_{2n} \\mid \\sigma\\tau = \\tau\\sigma\\}\n\t\t\t\t\\stackrel{\\mbox{\\tiny def}}{=} \\name{Norm}(\\tau),\n\t\t\t\\end{equation*}\n\t\t\tand the image of the subgroup $D_n \\subset B_n$ under $\\Phi$ is the\n\t\t\tintersection of $\\name{Norm}(\\tau)$ with group $S_{2n}^+ \n\t\t\t\\subset S_{2n}$ of even permutations. \n\t\t\\end{proposition}\n\t\t\n\t\tIntroduce a notation more convenient for our purposes: let $r_{ij}\n\t\t\\stackrel{\\mbox{\\tiny def}}{=} (i,j)(i+n,j+n)$ for all $1 \\le i \\ne j \\le 2n$, $i \\ne j\\pm n$;\n\t\taddition is modulo $2n$. So if $1 \\le i < j \\le n$ then $r_{ij} =\n\t\tr_{i+n,j+n} = \\Phi(s_{ij}^+)$ and $r_{i,j+n} = r_{i+n,j} =\n\t\t\\Phi(s_{ij}^-)$. Also denote $l_i \\stackrel{\\mbox{\\tiny def}}{=} (i,i+n)$ for all $1 \\le i \\le\n\t\t2n$; so, $l_i = l_{i+n} = \\Phi(\\ell_i)$ if $i \\le n$.\n\t\t\n\t\tA proof (rather elementary) of Proposition \\ref{Pp:Embed} is preceded\n\t\tby an explicit description of $\\name{Norm}(\\tau)$, which will be used\n\t\textensively throughout this article:\n\t\t\n\t\t\\begin{lemma}\\label{Lm:AlphaBeta}\n\t\t\tIf $x \\in \\name{Norm}(\\tau)$ then for every cycle $(\\gamma_1 \\dots \\gamma_m)$\n\t\t\tof its cycle decomposition one of the following is true:\n\t\t\t\n\t\t\t\\begin{enumerate}\n\t\t\t\t\\item\\label{It:Alpha} The cycle decomposition contains another\n\t\t\t\tcycle $(\\gamma_1' \\dots \\gamma_m')$ of the same length such that\n\t\t\t\t$\\gamma_i' = \\tau(\\gamma_i)$ for $i = 1 \\DT, s$.\n\t\t\t\t\\item\\label{It:Beta} $s = 2t$ is even and $\\gamma_{i+t} =\n\t\t\t\t\\tau(\\gamma_i)$ for all $i = 1 \\DT, t$.\n\t\t\t\\end{enumerate}\n\t\t\\end{lemma}\n\t\t\n\t\tIn the two cases we speak about {\\em $\\alpha$-pairs} of cycles and\n\t\tabout {\\em $\\beta$-cycles}, respectively.\n\t\t\n\t\t\\begin{proof}\n\t\t\tLet\n\t\t\t\\begin{equation*}\n\t\t\t\tx = (\\gamma_{11} \\dots \\gamma_{1m_1}) \\dots (\\gamma_{q1} \\dots\n\t\t\t\t\\gamma_{q m_q})\n\t\t\t\\end{equation*}\n\t\t\tbe the cycle decomposition. If $\\tau x=x\\tau$ then $\\tau x \\tau^{-1}\n\t\t\t= x$. So\n\t\t\t\\begin{equation*}\n\t\t\t\tx = (\\tau(\\gamma_{11}) \\dots \\tau(\\gamma_{1m_1})) \\dots\n\t\t\t\t(\\tau(\\gamma_{q1}) \\dots \\tau(\\gamma_{qm_q}))\n\t\t\t\\end{equation*}\n\t\t\tThe cycle decomposition is unique, so every cycle\n\t\t\t$(\\tau(\\gamma_{i1}) \\dots \\tau(\\gamma_{im_i}))$ must be equal to some\n\t\t\t$(\\gamma_{j_1} \\dots \\gamma_{j m_j})$. If these two cycles are\n\t\t\tdifferent, then they form an $\\alpha$-pair. If these\n\t\t\ttwo cycles are the same then there exists some $t$ such that\n\t\t\t$\\tau(\\gamma_k) = \\gamma_{k+t\\bmod m_i}$ for all $k = 1 \\DT, m_i$. Since\n\t\t\t$\\tau$ is an involution, one has $\\gamma_k = \\tau^2(\\gamma_i)\n\t\t\t=\\tau(\\gamma_{k+t \\bmod m_i}) = \\gamma_{k+2t \\bmod m_i}$ for all $k\n\t\t\t= 1 \\DT, m_i$, which implies $m_i = 2t$, and this is a $\\beta$-cycle.\n\t\t\\end{proof}\n\t\t\n\t\t\\begin{proof}[Proof of Proposition \\ref{Pp:Embed}]\n\t\t\tAny element $x \\in B_n$ is representable as $x = a_1\\dots a_N$ where\n\t\t\tevery $a_k$ is a reflection. Denote by $e_1 \\DT, e_n$ the standard\n\t\t\tbasis in $\\mathbb{R}^n$.\n\t\t\t\n\t\t\t\\begin{lemma}\\label{Lm:ActGen}\n\t\t\t\tFor any $k = 1 \\DT, n$ one has $x(e_k) = \\pm e_\\ell$ for some\n\t\t\t\t$\\ell = 1 \\DT, n$. If $x(e_k) = e_\\ell$ then $(\\tilde a_1 \\DT\n\t\t\t\t\\tilde a_N)(k) = \\ell$, and if $x(e_k) = -e_\\ell$ then\n\t\t\t\t$(\\tilde a_1 \\DT \\tilde a_N)(k) = \\t(\\ell) (= \\ell+n)$.\n\t\t\t\\end{lemma}\n\t\t\t\n\t\t\tThe lemma is proved by an immediate induction by $N$.\n\t\t\t\n\t\t\tLet us define now the map $\\Phi: B_n \\to S_{2n}$ as $\\Phi(x) =\n\t\t\t\\tilde a_1 \\dots \\tilde a_N$. If $x = b_1 \\dots b_M$ is another\n\t\t\trepresentation of $x$ as a product of reflections, then Lemma\n\t\t\t\\ref{Lm:ActGen} implies that $(\\tilde b_1 \\dots \\tilde b_M)(k) =\n\t\t\t(\\tilde a_1 \\dots \\tilde a_N)(k)$ for every $k$. Thus, the map\n\t\t\t$\\Phi$ is well-defined and is a group homomorphism.\n\t\t\t\n\t\t\tSuppose $\\Phi(x) = \\name{id}$ (that is, $x$ belongs to the kernel of\n\t\t\t$\\Phi$). Then $(\\tilde a_1 \\dots \\tilde a_N)(k) = k$ for every $k =\n\t\t\t1 \\DT, n$, and Lemma \\ref{Lm:ActGen} asserts then that $x(e_k) =\n\t\t\te_k$ for every $k$. Thus, $x = \\name{id}$, so $\\Phi$ is an\n\t\t\tembedding.\n\t\t\t\n\t\t\tShow now that $\\Phi(B_n) = \\name{Norm}(\\tau)$. Obviously, $\\Phi(B_n) \\subset\n\t\t\t\\name{Norm}(\\tau)$ because $\\Phi(s^\\pm_{ij})$ and $\\Phi(\\ell_i)$ commute\n\t\t\twith $\\tau$. Also, $\\Phi(D_n) \\subset S_{2n}^+$ because $r_{ij} \\in\n\t\t\tS_{2n}^+$.\n\t\t\t\n\t\t\tProve now that $\\name{Norm}(\\tau) \\subset \\Phi(B_n)$. Let $x \\in\n\t\t\t\\name{Norm}(\\tau)$. By Lemma \\ref{Lm:AlphaBeta}, $x$ is a product of\n\t\t\t$\\alpha$-pairs and of $\\beta$-cycles, so to prove that $x \\in\n\t\t\t\\Phi(B_n)$ it suffices to show that any $\\alpha$-pair and any\n\t\t\t$\\beta$-cycle are products of the elements $r_{ij}$ and\n\t\t\t$l_i$ ($1 \\le i < j\\le 2n$). This is the case:\n\t\t\t\\begin{equation*}\n\t\t\t\t(i_1 \\DT, i_m)(\\tau(i_1) \\DT, \\tau(i_m)) = r_{i_1 i_2}\n\t\t\t\t\\dots r_{i_{m-1} i_m}\n\t\t\t\\end{equation*}\n\t\t\tand\n\t\t\t\\begin{align*}\n\t\t\t\t(i_1 \\DT, i_m \\tau(i_1) \\DT, \\tau(i_m)) &= (i_1 \\DT,\n\t\t\t\ti_m)(\\tau(i_1) \\DT, \\tau(i_m)) (i_1\\tau(i_1))\\\\\n\t\t\t\t&= r_{i_1 i_2} \\dots r_{i_{m-1} i_m}\n\t\t\t\tl_{i_1}.\n\t\t\t\\end{align*}\n\t\t\tTo prove that $\\name{Norm}(\\tau) \\cap S_{2n}^+ \\subset \\Phi(D_n)$ observe\n\t\t\tthat the element $x \\in \\name{Norm}(\\tau)$ is even if and only it contains\n\t\t\tan even number of $\\beta$-cycles. (Indeed, an $\\alpha$-pair is always\n\t\t\tan even permutation, while a $\\beta$-cycle is odd because its length\n\t\t\tis even.) So it suffices to prove that the product of any two\n\t\t\t$\\beta$-cycles is a product of $r_{ij}$ (for\n\t\t\t$\\alpha$-cycles if was proved above). Indeed,\n\t\t\t\\begin{align*}\n\t\t\t\t(i_1 \\dots i_m \\tau(i_1) &\\dots \\tau(i_m)) (j_1 \\dots i_q\n\t\t\t\t\\tau(ij_1) \\dots \\tau(j_q))\\\\\n\t\t\t\t&= (i_1 \\dots i_m)(\\tau(i_1) \\dots \\tau(i_m)) (j_1 \\dots i_q)\n\t\t\t\t(\\tau(j_1) \\dots \\tau(j_q)) (i_1 \\tau(i_1))(j_1 \\tau(j_1).\n\t\t\t\\end{align*}\n\t\t\tThe product of the first four cycles is already proved to be a\n\t\t\tproduct of $r_{ij}$, while for the two remaining cycles\n\t\t\tone has\n\t\t\t\\begin{equation*}\n\t\t\t\t(i_1 \\tau(i_1))(j_1 \\tau(j_1)) = r_{ij} r_{i,\\tau(j)}.\n\t\t\t\\end{equation*}\n\t\t\\end{proof}\n\t\t\n\t\tTo save space, below we will not distinguish $B_n$ and $D_n$ from\n\t\ttheir images $\\Phi(B_n), \\Phi(D_n) \\subset S_{2n}$.\n\t\t\n\t\t\\subsection{Conjugacy classes}\\label{SSec:ConjClass}\n\t\t\n\t\tConjugacy classes in the permutation group $S_m$ are in one-to-one\n\t\tcorrespondence with partitions of $m$: two elements of $S_m$ are\n\t\tconjugate if and only if they have cycles of the same lengths\n\t\t(totalling $m$) in their cycle decomposition. A similar result for the\n\t\treflection groups $B_n$ and $D_n$ can be found in \\cite{Carter}. Quote\n\t\tit here for coompleteness; be warned that the original notation of\n\t\t\\cite{Carter} differ from the one used here.\n\t\t\n\t\tFix two partitions, $\\lambda = (\\lambda_1 \\DT, \\lambda_s)$ and\n\t\t$\\mu = (\\mu_1 \\DT, \\mu_t)$ such that $\\lmod \\lambda\\rmod + \\lmod\n\t\t\\mu\\rmod = \\lambda_1 \\DT+ \\lambda_s + \\mu_1 \\DT+ \\mu_t = n$, and\n\t\tconsider a set $C_{\\lambda\\mid\\mu}$ of elements $x \\in B_n \\subset\n\t\tS_{2n}$ such that their cycle decomposition contains\n\t\t\\begin{enumerate}\n\t\t\t\\item $\\alpha$-pairs of lengths (each cycle) $\\lambda_1, \\lambda_2\n\t\t\t\\DT, \\lambda_s$;\n\t\t\t\n\t\t\t\\item $\\beta$-cycles of lengths $2\\mu_1 \\DT, 2\\mu_t$ (recall that the\n\t\t\tlength should be even).\n\t\t\\end{enumerate}\n\t\t\n\t\t\\begin{proposition}[\\protect{\\cite[Proposition 25]{Carter}}]\\label{Pp:ConjClass}\n\t\t\tThe set $C_{\\lambda\\mid\\mu} \\subset B_n$ is a conjugacy class. Every\n\t\t\tconjugacy class in $B_n$ is $C_{\\lambda\\mid\\mu}$ for some $\\lambda$ and\n\t\t\t$\\mu$ such that $\\lmod\\lambda\\rmod + \\lmod\\mu\\rmod = n$.\n\t\t\\end{proposition}\n\t\t\n\t\tFor $D_n$ the answer is slightly more complicated. Take a partition\n\t\t$\\lambda = (\\lambda_1 \\DT, \\lambda_s)$ such that $\\lmod\\lambda\\rmod =\n\t\tn$ and all $\\lambda_i$ are even. For an element $\\sigma \\in\n\t\tC_{\\lambda\\mid\\emptyset}$ write its cycle decomposition\n\t\t\\begin{equation}\\label{Eq:2LambdaEmpty}\n\t\t\t\\sigma = (\\gamma_{11} \\dots \\gamma_{1\\lambda_1})(\\tau(\\gamma_{11})\n\t\t\t\\dots \\tau(\\gamma_{1\\lambda_1})) \\dots (\\gamma_{s1} \\dots \\gamma_{s\\lambda_s})(\\tau(\\gamma_{s1})\n\t\t\t\\dots \\tau(\\gamma_{s\\lambda_s}))\n\t\t\\end{equation}\n\t\tsatisfying two conditions: cycles forming an $\\alpha$-pair always\n\t\tstand side by side, and their matching elements $\\gamma_{ij} $ and\n\t\t$\\tau(\\gamma_{ij} ) $ occupy the same positions in them. Now consider\n\t\ta permutation $\\gamma = (\\gamma_{11} \\dots \\gamma_{s\\lambda_s} ) \\in\n\t\tS_{2n}$ where the numbers $\\gamma_{ij}$ are written exactly in the\n\t\tsame order as in \\eqref{Eq:2LambdaEmpty}. We write $\\sigma \\in\n\t\tC_{\\lambda\\mid\\emptyset}^+$ if $\\gamma$ is even, and $\\sigma \\in\n\t\tC_{\\lambda\\mid\\emptyset}^-$ if $\\gamma$ is odd. Representation of\n\t\t$\\sigma$ in the form \\eqref{Eq:2LambdaEmpty} is not unique but, as it\n\t\tis easy to see, the parity of $\\gamma$ does not depend on a particular\n\t\tchoice. (Recall, all the cycles in \\eqref{Eq:2LambdaEmpty} have even\n\t\tlength.)\n\t\t\n\t\t\\begin{proposition}[\\protect{\\cite[Proposition 25]{Carter}}]\\label{Pp:ConjD}\\strut\n\t\t\t\\begin{enumerate}\n\t\t\t\t\\item If the partition $\\mu$ contains an even number of parts then\n\t\t\t\tthe conjugacy class $C_{\\lambda\\mid\\mu} \\subset B_n$ lies in $D_n$;\n\t\t\t\tif the number of parts is odd then $C_{\\lambda\\mid\\mu}$ does not\n\t\t\t\tintersect $D_n$.\n\t\t\t\t\n\t\t\t\t\\item If $\\mu \\ne \\emptyset$ and the number of parts of $\\mu$ is\n\t\t\t\teven then $C_{\\lambda\\mid\\mu}$ is a conjugacy class in $D_n$.\n\t\t\t\t\n\t\t\t\t\\item If $\\lambda$ is a partition of $n$ containing at least one odd\n\t\t\t\tpart $\\lambda_i$ then $C_{\\lambda\\mid\\emptyset}$ is a conjugacy class\n\t\t\t\tin $D_n$.\n\t\t\t\t\n\t\t\t\t\\item If $\\lambda$ is a partition of $n$ such that all its parts are\n\t\t\t\teven then $C_{\\lambda \\mid \\emptyset}$ splits into two conjugacy\n\t\t\t\tclasses in $D_n$, $C_{\\lambda\\mid\\emptyset}^+$ and\n\t\t\t\t$C_{\\lambda\\mid\\emptyset}^-$. \n\t\t\t\\end{enumerate}\n\t\t\tAny conjugacy class in $D_n$ is one of the classes listed above.\n\t\t\\end{proposition}\n\t\t\n\t\t\\begin{corollary}\\label{Cr:ConjOdd}\n\t\t\t\\begin{enumerate}\n\t\t\t\t\\item Let $\\lambda, \\mu$ be partitions such that $\\lmod\\lambda\\rmod\n\t\t\t\t+ \\lmod\\mu\\rmod = n$, $\\#\\mu$ is even and either $\\mu \\ne\n\t\t\t\t\\emptyset$ or at least one of he parts of $\\lambda$ is odd. Then\n\t\t\t\tfor any $x \\in B_n$ and any $\\sigma \\in C_{\\lambda\\mid\\mu} \\subset\n\t\t\t\tD_n$ one has $x\\sigma x^{-1} \\in C_{\\lambda\\mid\\mu}$. \n\t\t\t\t\\item Let $\\lambda$ be a partition of $n$ such that all its parts\n\t\t\t\tare even, and let $x \\in B_n$. Then for any $\\sigma \\in\n\t\t\t\tC_{\\lambda\\mid\\emptyset}^+$ one has $x \\sigma x^{-1} \\in\n\t\t\t\tC_{\\lambda\\mid\\emptyset}^+$ if $x \\in D_n \\subset B_n$ (that is,\n\t\t\t\t$x \\in S_{2n}$ is an even permutation) and $x \\sigma x^{-1} \\in\n\t\t\t\tC_{\\lambda\\mid\\emptyset}^-$ otherwise; the picture for $\\sigma \\in\n\t\t\t\tC_{\\lambda\\mid\\emptyset}^-$ is symmetric.\n\t\t\t\\end{enumerate}\n\t\t\\end{corollary}\n\t\t\n\t\tIn particular, $C_{{1^{n-2}2^1}\\mid\\varnothing} \\subset D_n \\subset\n\t\tB_n$ consists of all the reflections $r_{ij}$, and $C_{{1^{n-1}}\\mid\n\t\t\t1}\\subset B_n$, of all the reflections $l_i$.\n\t\t\n\t\t\\section{Hurwitz numbers}\\label{Sec:HurwNum}\n\t\t\n\t\t\\subsection{Definitions}\n\t\t\n\t\tFix a pair of partitions $\\lambda, \\mu$ with $\\lmod\\lambda\\rmod +\n\t\t\\lmod\\mu\\rmod = n$. Let $C_{\\lambda\\mid\\mu} \\subset B_n$ be the conjugacy\n\t\tclass defined above.\n\t\t\n\t\t\\begin{definition}\\label{Df:HurwNumb}\n\t\t\tA sequence of reflections $(\\sigma_1 \\DT, \\sigma_{m+\\ell})$ of the\n\t\t\tgroup $B_n$ is said to have {\\em profile} $(\\lambda,\\mu,m,\\ell)$ if\n\t\t\t$\\#\\{p \\mid \\sigma_p = r_{ij}, 1 \\le i < j \\le 2n\\} = m$,\n\t\t\t$\\#\\{p \\mid \\sigma_p = l_i, 1 \\le i \\le 2n\\} = \\ell$ and\n\t\t\t$\\sigma_1 \\dots \\sigma_{m+\\ell} \\in C_{\\lambda \\mid \\mu}$. The {\\em\n\t\t\t\tHurwitz numbers for the group $B_n$} are $h_{m,\\ell,\\lambda,\\mu}^B\n\t\t\t= \\frac{1}{n!} \\#\\{(\\sigma_1 \\DT, \\sigma_{m+\\ell}) \\text{ is a\n\t\t\t\tsequence of profile } (\\lambda,\\mu,m,\\ell)\\}$.\n\t\t\\end{definition}\n\t\t\n\t\tFor the group $D_n$ we use the same numbers, in case they make sense:\n\t\t\n\t\t\\begin{definition}\\label{Df:D-Hurwitz}\n\t\t\tLet $m$ be a positive integer, and $\\lambda$ and $\\mu$, partitions\n\t\t\twhere the number of parts $\\#\\mu$ is even. The Hurwitz number for\n\t\t\tthe group $D_n$ is defined as $h_{m,\\lambda,\\mu}^D =\n\t\t\th_{m,0,\\lambda,\\mu}^B$.\n\t\t\\end{definition}\n\t\t\n\t\tIn other words, $h_{m,\\lambda,\\mu}^D = \\frac{1}{n!} \\#\\{(\\sigma_1 \\DT,\n\t\t\\sigma_m) \\text{ has profile } (\\lambda,\\mu,m)\\}$ where $\\sigma_p$ are\n\t\treflections in the group $D_n$ (that is, $\\sigma_p = r_{ij}$ for some\n\t\t$i$ and $j$), and the profile means that $\\sigma_1 \\dots \\sigma_m \\in\n\t\tC_{\\lambda\\mid\\mu} \\subset D_n$.\n\t\t\n\t\t\\begin{remark}\\label{Rm:DHurwRefine}\n\t\t\tRecall (Proposition \\ref{Pp:ConjD}) that if $\\mu = \\emptyset$ and\n\t\t\tall the parts of the partition $\\lambda$ are even then\n\t\t\t$C_{\\lambda\\mid\\emptyset}$ splits into two conjugacy classes,\n\t\t\t$C_{\\lambda\\mid\\emptyset}^+$ and $C_{\\lambda\\mid\\emptyset}^-$.\n\t\t\t\n\t\t\tDenote by $\\Sigma_{m,\\lambda}^+$ and $\\Sigma_{m,\\lambda}^-$ the sets\n\t\t\tof $m$-tuples of reflections $\\sigma_1 \\DT, \\sigma_m \\in D_n$ such\n\t\t\tthat $\\sigma_1 \\dots \\sigma_m \\in C_{\\lambda\\mid\\emptyset}^+$\n\t\t\t(resp., $C_{\\lambda\\mid\\emptyset}^-$). One can denote\n\t\t\t$h_{m,\\lambda,\\emptyset\\pm}^D \\stackrel{\\mbox{\\tiny def}}{=} \\frac{1}{n!}\n\t\t\t\\#\\Sigma_{m,\\lambda}^\\pm$, so that $h_{m,\\lambda,\\emptyset+}^D +\n\t\t\th_{m,\\lambda,\\emptyset-}^D = h_{m,\\lambda,\\emptyset}^D$. It follows\n\t\t\tfrom Corollary \\ref{Cr:ConjOdd}, though, that if $x \\in B_n$ but $x\n\t\t\t\\notin D_n$ then the map sending an $m$-tuple $(\\sigma_1 \\DT,\n\t\t\t\\sigma_m)$ to $(x\\sigma_1x^{-1} \\DT, x\\sigma_mx^{-1})$ is a\n\t\t\tbijection between $\\Sigma_{m,\\lambda}^+$ and\n\t\t\t$\\Sigma_{m,\\lambda}^-$. Thus, $h_{m,\\lambda,\\emptyset+}^D =\n\t\t\th_{m,\\lambda,\\emptyset-}^D = \\frac12 h_{m,\\lambda,\\emptyset}^D$, so\n\t\t\tconsidering $h_{m,\\lambda,\\emptyset\\pm}^D$ makes little sense.\n\t\t\\end{remark}\n\t\t\n\t\tUp to the end of this section we consider the Hurwitz numbers for the\n\t\tgroup $B_n$ only.\n\t\t\n\t\tDenote by $\\CSum{\\lambda}{\\mu} \\stackrel{\\mbox{\\tiny def}}{=} \\frac{1}{\\#C_{\\lambda\\mid\\mu}}\n\t\t\\sum_{x \\in C_{\\lambda\\mid\\mu}} x \\in \\mathbb{C}[B_n]$ the normalized\n\t\tclass sum. One has $\\CSum{\\lambda}{\\mu} \\in Z[B_n]$ (the\n\t\tcenter of the group algebra of $B_n$); by Proposition\n\t\t\\ref{Pp:ConjClass}, $\\CSum{\\lambda}{\\mu}$ form a basis in\n\t\t$Z[B_n]$. Consider now a ring of polynomials $\\mathbb{C}[p,q]$\n\t\twhere $p = (p_1, p_2, \\dots)$ and $q = (q_1, q_2, \\dots)$ are two\n\t\tinfinite sets of variables. The ring is graded by the total degree\n\t\twhere one assumes $\\deg p_k=\\deg q_k =k$ for all $k=1,2,\\dots$. The\n\t\tmap $\\Psi$ defined by\n\t\t\\begin{equation}\\label{Eq:DefIsoB}\n\t\t\t\\Psi(\\CSum{\\lambda}{\\mu}) = p_\\lambda q_\\mu \\stackrel{\\mbox{\\tiny def}}{=} p_{\\lambda_1}\n\t\t\t\\dots p_{\\lambda_s} q_{\\mu_1} \\dots q_{\\mu_t}\n\t\t\\end{equation}\n\t\testablishes an isomorphism between $Z[B_n]$ and the\n\t\thomogeneous component $\\mathbb{C}[p,q]_n$ of total degree $n$.\n\t\t\n\t\tNow denote by\n\t\t\\begin{equation}\\label{Eq:DefT1}\n\t\t\t\\mathcal T_1 := \\frac12 \\sum_{1 \\le i < j \\le 2n}\n\t\t\tr_{ij} = \\#C_{2 1^{n-2}\\mid\\emptyset} \\cdot \\CSum{1^{n-2}2}{\\emptyset}\n\t\t\\end{equation}\n\t\tand\n\t\t\\begin{equation*}\n\t\t\t\\mathcal T_2 := \\frac12 \\sum_{1 \\le i \\le 2n} l_i = \\#C_{1^{n-1}\\mid1}\n\t\t\t\\cdot \\CSum{1^{n-1}}{1}\n\t\t\\end{equation*}\n\t\tsums of all elements of the conjugacy classes containing reflections\n\t\t(recall that $r_{ij} = r_{i+n,j+n}$, and $l_i = l_{i+n}$, so every\n\t\telement is repeated twice in these sums; hence the factor\n\t\t$\\frac12$). The elements $\\mathcal T_1$ and $\\mathcal T_2$ belong to\n\t\t$Z[B_n]$, so one can consider linear operators $T_1, T_2: Z[B_n] \\to\n\t\tZ[B_n]$ of multiplication by $\\mathcal T_1$ and $\\mathcal T_2$,\n\t\trespectively. Obviously, $T_1$ and $T_2$ commute.\n\t\t\n\t\tConsider now linear operators $\\mathcal{CJ}_1, \\mathcal{CJ}_2: \\mathbb{C}[p,q]_n \\to\n\t\t\\mathbb{C}[p,q]_n$ making the following diagrams commutative:\n\t\t\n\t\t\\begin{equation}\\label{Eq:CJViaGAlg}\n\t\t\t\\xymatrix{\n\t\t\t\tZ\\mathbb{C}[B_n] \\ar[r]^{\\times T_i} \\ar[d]^\\simeq &\n\t\t\t\tZ\\mathbb{C}[B_n] \\ar[d]^\\simeq \\\\\n\t\t\t\t\\mathbb{C}[p,q]_n \\ar[r]_{\\mathcal{CJ}_i} &\n\t\t\t\t\\mathbb{C}[p,q]_n \n\t\t\t}, \\qquad i = 1,2.\n\t\t\\end{equation}\n\t\t\n\t\tLet now $\\lambda, \\lambda', \\mu, \\mu'$ be partitions such that\n\t\t$\\lmod\\lambda\\rmod + \\lmod\\mu\\rmod = \\lmod\\lambda'\\rmod +\n\t\t\\lmod\\mu'\\rmod = n$. Take an element $\\sigma_* \\in C_{\\lambda\\mid\\mu}$\n\t\tand define the {\\em multiplicity} $\\langle \\lambda, \\mu \\mid \\lambda',\n\t\t\\mu'\\rangle_1$ as the number of reflections $u \\in C_{2^11^{n-2} \\mid\n\t\t\t\\varnothing}$ (that is, $u = r_{ij}$) such that $u\n\t\t\\sigma_* \\in C_{\\lambda'\\mid\\mu'}$; the multiplicity $\\langle \\lambda,\n\t\t\\mu \\mid \\lambda', \\mu'\\rangle_2$ is defined in the same way with\n\t\t$u\\in C_{1^{n-1}\\mid1}$ (that is, $u = l_i$) instead.\n\t\t\n\t\t\\begin{lemma}\\label{Lm:Conj}\n\t\t\tMultiplicities do not depend on the choice of $\\sigma_*$.\n\t\t\\end{lemma}\n\t\t\n\t\t\\begin{proof}\n\t\t\tDenote by $S_1(\\sigma_*;\\lambda',\\mu') \\subset\n\t\t\tC_{1^{n-2}2^1\\mid\\emptyset}$ the set of reflections $u$ such that\n\t\t\t$u\\sigma_* \\in C_{\\lambda'\\mid\\mu'}$. Take $\\sigma' = x\\sigma_*\n\t\t\tx^{-1} \\in C_{\\lambda\\mid\\mu}$ where $x \\in B_n$. If $u \\in\n\t\t\tS(\\sigma';\\lambda',\\mu')$ then $u \\sigma' = u x\\sigma_* x^{-1} = x\n\t\t\t((x^{-1}ux) \\sigma_*) x^{-1} \\in C_{\\lambda'\\mid\\mu'}$, which is\n\t\t\tequivalent to $(x^{-1}ux) \\sigma_* \\in C_{\\lambda'\\mid\\mu'}$, that\n\t\t\tis, to $x^{-1}ux \\in S(\\sigma_*;\\lambda',\\mu')$. Thus, conjugation\n\t\t\t$u \\mapsto x^{-1}ux$ is a one-to-one map sending\n\t\t\t$S_1(\\sigma';\\lambda',\\mu')$ to $S_1(\\sigma_*;\\lambda',\\mu')$, and\n\t\t\ttherefore these two sets contain the same number of elements\n\t\t\t$\\langle\\lambda,\\mu \\mid \\lambda',\\mu'\\rangle_1$. The reasoning for\n\t\t\tthe multiplicity $\\langle\\lambda,\\mu \\mid \\lambda',\\mu'\\rangle_2$ is\n\t\t\tthe same.\n\t\t\\end{proof}\n\t\t\n\t\t\\begin{theorem}\\label{Th:TViaMult}\n\t\t\t$T_i\\CSum{\\lambda}{\\mu} = \\sum_{\\lambda',\\mu'} \\langle \\lambda,\\mu\n\t\t\t\\mid \\lambda',\\mu'\\rangle_i \\CSum{\\lambda'}{\\mu'}$ for $i = 1,2$.\n\t\t\\end{theorem}\n\t\t\n\t\t\\begin{proof}\n\t\t\t\\begin{equation}\\label{Eq:ActionT}\n\t\t\t\tT_1 \\CSum{\\lambda}{\\mu} = \\frac{1}{\\#C_{\\lambda\\mid\\mu}}\n\t\t\t\t\\sum_{\\sigma \\in C_{\\lambda\\mid\\mu}} T_1\\sigma =\n\t\t\t\t\\frac{1}{\\#C_{\\lambda\\mid\\mu}} \\sum_{\\sigma \\in\n\t\t\t\t\tC_{\\lambda\\mid\\mu}}\\sum_{u \\in C_{1^{n-2}2^1\\mid\\varnothing}} u\n\t\t\t\t\\sigma .\n\t\t\t\\end{equation}\n\t\t\tIt follows from Lemma \\ref{Lm:Conj} that \\eqref{Eq:ActionT} is an\n\t\t\tarithmetical mean of identical summands, so it is equal to each of\n\t\t\tthem:\n\t\t\t\\begin{equation*}\n\t\t\t\tT_1 \\CSum{\\lambda}{\\mu} = \\sum_{\\lambda',\\mu'} \\sum_{\\substack\n\t\t\t\t\t{u \\in C_{1^{n-2}2^1 \\mid \\varnothing}\\\\\n\t\t\t\t\t\tu\\sigma_* \\in C_{\\lambda'\\mid\\mu'}}}\n\t\t\t\tu \\sigma_* \n\t\t\t\\end{equation*}\n\t\t\tfor any fixed $\\sigma_* \\in C_{\\lambda\\mid\\mu}$. Using Lemma\n\t\t\t\\ref{Lm:Conj} again, one obtains\n\t\t\t\\begin{align*}\n\t\t\t\tT_1 \\CSum{\\lambda}{\\mu} &= \\sum_{\\lambda',\\mu'} \\sum_{\\substack{\n\t\t\t\t\t\tu \\in C_{1^{n-2}2^1 \\mid \\varnothing}\\\\\n\t\t\t\t\t\tu\\sigma_* \\in C_{\\lambda'\\mid\\mu'}}}\n\t\t\t\t\\frac{1}{\\#C_{\\lambda'\\mid\\mu'}}\n\t\t\t\t\\sum_{\\tau \\in C_{\\lambda'\\mid\\mu'}} \\tau\\\\\n\t\t\t\t&= \\sum_{\\lambda',\\mu'} \\#(\\{u \\in C_{1^{n-2}2^1 \\mid\\varnothing}:\n\t\t\t\tu\\sigma_* \\in C_{\\lambda'\\mid\\mu'}\\} \\CSum{\\lambda'}{\\mu'}\\\\\n\t\t\t\t&= \\sum_{\\lambda',\\mu'} \\langle \\lambda,\\mu \\mid\n\t\t\t\t\\lambda',\\mu'\\rangle \\CSum{\\lambda'}{\\mu'}.\n\t\t\t\\end{align*}\n\t\t\tThe proof for $T_2$ is the same with $C_{1^{n-1}\\mid 1}$ instead of\n\t\t\t$C_{1^{n-2}2^1\\mid \\varnothing}$.\n\t\t\\end{proof}\n\t\t\n\t\t\\subsection{Generating function}\n\t\t\n\t\tConsider the following generating function for Hurwitz numbers of the\n\t\tgroup $B_n$:\n\t\t\\begin{equation*}\n\t\t\t\\mathcal{H}^B(\\beta,\\gamma,p,q) = \\sum_{m,\\ell}\\sum_{\\lambda,\\mu}\n\t\t\t\\frac{ h_{m,\\ell,\\lambda,\\mu}^B}{m!\\ell!}\n\t\t\tp_{\\lambda}q_{\\mu}\\beta^m\\gamma^\\ell.\n\t\t\\end{equation*}\n\t\t\n\t\t\\begin{theorem}\\label{Th:CJequ}\n\t\t\t$\\mathcal{H}^B$ satisfies the {\\em cut-and-join equations}\n\t\t\t\\begin{equation}\\label{Eq:CJ}\n\t\t\t\t\\pder{\\mathcal{H}^B}{\\beta} = \\mathcal{CJ}_1(\\mathcal{H}^B) \\quad \\text{and}\n\t\t\t\t\\quad \\pder{\\mathcal{H}^B}{\\gamma} = \\mathcal{CJ}_2(\\mathcal{H}^B)\n\t\t\t\\end{equation}\n\t\t\\end{theorem}\n\t\t\n\t\t\\begin{proof}\n\t\t\tFix a positive integer $n$ and denote by $\\mathcal{H}_n$ a degree\n\t\t\t$n$ homogeneous component of $\\mathcal{H}^B$. The cut-and-join\n\t\t\toperators preserve the degree, so $\\mathcal{H}^B$ satisfies the\n\t\t\tcut-and-join equations if and only if $\\mathcal{H}_n$ does (for each\n\t\t\t$n$).\n\t\t\t\n\t\t\tLet $i$ be $1$ or $2$, and let\n\t\t\t\\begin{equation*}\n\t\t\t\t\\mathcal{G}_n \\stackrel{\\mbox{\\tiny def}}{=} \\sum_{m,\\ell \\ge 0}\n\t\t\t\t\\sum_{\\lambda,\\mu:\\mid\\lambda\\mid +\\mid\\mu\\mid = n}\n\t\t\t\t\\frac{n!h_{m,\\ell,\\lambda,\\mu}^{B}}{m!\\ell!} \\CSum{\\lambda}{\\mu}\n\t\t\t\t\\beta^m \\gamma^\\ell \\in\\mathbb{C}[B_n]\n\t\t\t\\end{equation*}\n\t\t\tAn elementary combinatorial reasoning gives\n\t\t\t\\begin{equation*}\n\t\t\t\t\\mathcal{G}_n = \\sum_{m\\geq 0}\n\t\t\t\t\\frac{\\beta^m\\gamma^\\ell}{m!\\ell!}T_1^m T_2^\\ell(e_n)\n\t\t\t\\end{equation*}\n\t\t\twhere $e_n \\in B_n$ is the unit element. Clearly\n\t\t\t\\begin{equation*}\n\t\t\t\tT_1(\\mathcal{G}_n) = \\sum_{m\\geq 0}\n\t\t\t\t\\frac{\\beta^m\\gamma^\\ell}{m!\\ell!}(T_1)^{m+1}(e_{n}) = \\sum_{m \\ge\n\t\t\t\t\t1} \\frac{\\beta^{m-1}\\gamma^\\ell}{(m-1)!\\ell!} (T_1)^{m}(e_{n}) =\n\t\t\t\t\\pder{\\mathcal{G}_n}{\\beta}.\n\t\t\t\\end{equation*}\n\t\t\tand, similarly, $T_1(\\mathcal{G}_n) =\n\t\t\t\\pder{\\mathcal{G}_n}{\\gamma}$. Applying the isomorphism $\\Psi$ one\n\t\t\tobtains $\\Psi T_1(\\mathcal{G}_n) = \\Psi(\\pder{\\mathcal{G}_n}{\\beta})\n\t\t\t= \\pder{}{\\beta} \\Psi(\\mathcal{G}_n)$. One has $\\Psi(\\mathcal{G}_n)\n\t\t\t= \\mathcal{H}_n$, hence $\\pder{}{\\beta} \\Psi(\\mathcal{G}_n) =\n\t\t\t\\pder{\\mathcal{H}_n}{\\beta}$; similarly, $\\pder{}{\\gamma}\n\t\t\t\\Psi(\\mathcal{G}_n) = \\pder{\\mathcal{H}_n}{\\gamma}$. By the\n\t\t\tdefinition of the cut-and-join operators, $\\Psi T_1(\\mathcal{G}_n) =\n\t\t\t\\mathcal{CJ}_1(\\Psi(\\mathcal{G}_n)) = \\mathcal{CJ}_1({\\mathcal{H}_n})$ and the same\n\t\t\tfor $T_2$ and $\\mathcal{CJ}_2$. Equalities \\eqref{Eq:CJ} follow.\n\t\t\\end{proof}\n\t\t\n\t\t\\begin{corollary}\\label{Cr:GenFun}\n\t\t\t\\begin{equation}\\label{Eq:GenFunc}\n\t\t\t\t\\mathcal{H}^B(\\beta,\\gamma,p,q) =\n\t\t\t\te^{\\beta\\mathcal{CJ}_1+\\gamma\\mathcal{CJ}_2}e^{p_1}\n\t\t\t\\end{equation}\n\t\t\\end{corollary}\n\t\t\n\t\t\\begin{proof}\n\t\t\tIt follows from Definition \\ref{Df:HurwNumb} that\n\t\t\t$h^B_{0,0,\\lambda,\\mu} = \\frac{1}{n!}$ if $\\lambda = 1^n$ and\n\t\t\t$\\mu=\\emptyset$, and $h^B_{0,0,\\lambda,\\mu} = 0$ otherwise. Thus\n\t\t\t$\\mathcal{H}^B(0,0,p,q) = e^{p_1}$ and \\eqref{Eq:GenFunc} follows\n\t\t\tfrom Theorem \\ref{Th:CJequ}.\n\t\t\\end{proof}\n\t\t\n\t\t\\section{Explicit formulas}\\label{Sec:Expl}\n\t\t\n\t\t\\subsection{Multiplication of an element $x \\in B_n \\subset\n\t\t\tS_{2n}$ by a reflection}\\label{SSec:CycleStruct}\n\t\t\n\t\tLet $\\omega \\in S_n$ and $1 \\le a < b \\le n$. The cycle structure of\n\t\tthe product $\\omega' = (ab)\\omega$ depend on the cyclic structure of\n\t\t$\\omega$ and positions of $a$ and $b$ as follows. If $a$ and $b$\n\t\tbelong to the same cycle $(a \\DT, u, b \\DT, v)$ of $\\omega$, then the\n\t\tcycle splits in $\\omega'$ into two: $(a \\DT, u)$ and $(b \\DT, v)$ (``a\n\t\tcut''). If $a$ and $b$ belong to different cycles then the opposite\n\t\tthing happens: they glue together in $\\omega'$ (``a join''). The\n\t\tcycles containing neither $a$ nor $b$ are the same for $\\omega$ and\n\t\t$\\omega'$.\n\t\t\n\t\tLet now $\\omega \\in B_n \\subset S_{2n}$. The cyclic structure of $\\omega'\n\t\t= \\sigma\\omega$ where $\\sigma$ is a reflection (i.e.\\ $\\sigma =\n\t\tr_{ab} = (ab)(a+n,b+n)$ (addition modulo $n$) for $1 \\le a\n\t\t< b \\le 2n$, or $\\sigma = l_a = (a,a+n)$ for $1\\leq\n\t\ta\\leq 2n$) depends on the cyclic structure of $\\omega$ ($\\alpha$-pairs\n\t\tand $\\beta$-cycles) and the position of the points $a, b$ as shown in\n\t\tTable \\ref{Fig:Mult} below.\n\t\t\n\t\tIn this table, hexagons represent $\\beta$-cycles and the pairs of\n\t\ttriangles are $\\alpha$-pairs; numbers inside denote the lengths of the\n\t\tcycles. The two-headed dashed arrows show the position of $a$ and $b$\n\t\tof the reflection $\\sigma$. The solid arrows join cyclic structures of\n\t\t$\\omega$ and $\\sigma\\omega$ for $\\sigma = r_{ab}$, and the empty\n\t\tarrows, for $\\sigma = l_a$ (recall that $\\sigma$ is an involution,\n\t\tso the arrows are two-headed). For example, the first diagram shows\n\t\tthe multiplication by $(a,b)(\\t(a),\\t(b))$ of a pair of $\\beta$-cycles\n\t\tcontaining $a$ and $b$, respectively.\n\t\t\n\t\tThe boxed number on the bottom left corner is the multiplicity\n\t\t$\\langle\\lambda,\\mu \\mid \\lambda',\\mu'\\rangle$ between conjugacy\n\t\tclasses containing $\\omega$ and $\\sigma\\omega$.\n\t\t\n\t\t\n\t\t\\begin{figure} \n\t\t\t\\includegraphics[scale=.6]{ActionFinale.pdf} \n\t\t\t\\includegraphics[scale=.6]{ActionFinale2.pdf}\n\t\t\t\\caption{Multiplication of cycles by a reflection}\\label{Fig:Mult}\n\t\t\\end{figure}\n\t\t\n\t\t\\subsection{Multiplicities}\\label{SSec:Mult}\n\t\t\n\t\tLet $\\omega \\in C_{\\lambda\\mid\\mu}$ where $\\lambda = 1^{c_1} \\dots n^{c_n}$\n\t\tand $\\mu = 1^{d_1} \\dots n^{d_n}$ --- that is, for every $k = 1 \\DT,\n\t\tn$ the element $u$ contains $c_k$ $\\alpha$-pairs of length $k$ and\n\t\t$d_k$ separate $\\beta$-cycles of length $2k$. Calculate, for all\n\t\t$\\lambda' = 1^{c_1'} \\dots n^{c_n'}$ and $\\mu = 1^{d_1'} \\dots\n\t\tn^{d_n'}$, the multiplicities $\\langle \\lambda',\\mu' \\mid \\lambda,\n\t\t\\mu\\rangle_1$ and $\\langle \\lambda',\\mu' \\mid \\lambda,\n\t\t\\mu\\rangle_2$. In Section \\ref{SSec:CycleStruct} we listed cases when\n\t\tthis multiplicities may be nonzero: these are Cases\n\t\t\\ref{It:OneBeta}--\\ref{It:2join} below for $\\langle\\;\\mid\\;\\rangle_1$ and\n\t\tCases \\ref{It:LCut} and \\ref{It:LJoin} for $\\langle\\;\\mid\\;\\rangle_2$.\n\t\t\n\t\t{\\def Case \\theenumi. {Case \\theenumi.}\n\t\t\t\n\t\t\t\\begin{enumerate} \n\t\t\t\t\n\t\t\t\t\\item\\label{It:OneBeta} The number of possible positions for $a$ is\n\t\t\t\tthe total number of elements in all the $\\beta$-cycles of length\n\t\t\t\t$2(i+j)$, that is, $2(i+j) d_{i+j}$. Since one knows the\n\t\t\t\tlengths of $\\alpha$-pairs in $\\sigma \\omega$, the position of $b$ is\n\t\t\t\tunique, once the position of $a$ is chosen; the same is for $\\t(a)$\n\t\t\t\tand $\\t(b)$. Doing like this, one counts every reflection $\\sigma =\n\t\t\t\tr_{ab}$ twice, so that the multiplicity is $(i+j)d_{i+j}$.\n\t\t\t\t\n\t\t\t\t\\item Again there are $2i c_i$ possible positions for $a$ and $2j d_j$\n\t\t\t\tpossible positions for $b$; we are to divide by $2$ again by the\n\t\t\t\tsame reason and the multiplicity is $2ijc_i d_j$.\n\t\t\t\t\n\t\t\t\t\\item Like in Case \\ref{It:OneBeta}, the multiplicity is $2ij d_i\n\t\t\t\td_j$. .\n\t\t\t\t\n\t\t\t\t\\item\\label{It:MatchAlphas} The number of possible positions for $a$\n\t\t\t\tis $2(i+j) c_{i+j}$; the position of $b$ is uniquely determined by\n\t\t\t\tthe position of $a$ like in Case \\ref{It:OneBeta}. Again, divide by\n\t\t\t\t$2$, to obtain the multiplicity $(i+j) c_{i+j}$.\n\t\t\t\t\n\t\t\t\t\\item Same as case \\ref{It:MatchAlphas}, the multiplicity is\n\t\t\t\t$(i+j) c_{i+j}$.\n\t\t\t\t\n\t\t\t\t\\item\\label{It:2join} The number of possible positions for $a$ is the\n\t\t\t\ttotal number of the elements of all the $\\alpha$-pairs of length\n\t\t\t\t$i$, that is, $2i c_i$. Similarly, the number of possible positions\n\t\t\t\tfor $b$ is $2j c_j$, so the multiplicity is $2ij c_i c_j$.\n\t\t\t\t\n\t\t\t\t\\item\\label{It:LCut} The number of possible positions for $a$ is\n\t\t\t\t$2i d_i$, for the same reason as above, we have to divide by $2$,\n\t\t\t\tso that the multiplicity is $i d_i$\n\t\t\t\t\n\t\t\t\t\\item\\label{It:LJoin} The number of possible positions for $a$ is\n\t\t\t\t$2i c_i$, as in the previous cases, we have to divide by $2$, so that\n\t\t\t\tthe multiplicity is $i c_i$\n\t\t\t\\end{enumerate}\n\t\t}\n\t\t\n\t\tWrite now the terms of $\\mathcal{CJ}_1$ and $\\mathcal{CJ}_2$ explicitly. It follows from\n\t\tTheorem \\ref{Th:TViaMult} that $\\mathcal{CJ}_i p_\\lambda q_\\mu =\n\t\t\\sum_{\\lambda',\\mu'} \\langle\\lambda,\\mu\\mid\\lambda',\\mu'\\rangle_i\n\t\tp_{\\lambda'} q_{\\mu'}$ for $i\\in \\{1,2\\}$.\n\t\t\n\t\tLet $\\lambda',\\mu'$ be as at diagram \\ref{It:OneBeta} in Figure\n\t\t\\ref{Fig:Mult}. The monomial $p_\\lambda q_\\mu$ contains\n\t\t$q_{i+j}^{d_{i+j}}$; The exponent at $q_{i+j}$ in the monomial\n\t\t$p_{\\lambda'} q_{\\mu'}$ is less by $1$, and the exponents at $p_i$ and\n\t\t$q_j$ are greater by $1$. We will have $q_{i+j} \\pder{}{p_i}\n\t\t\\pder{}{q_j}$ for this; to get a correct multiplicity, put the\n\t\tcoefficient $(i+j)$ before.\n\t\t\n\t\tSimilar reasoning for the remaining cases gives:\n\t\t\n\t\t\\begin{multline}\\label{Eq:CJ_1}\n\t\t\t\\mathcal{CJ}_1 = \\sum_{i,j=1}^\\infty \\biggl( (i+j) p_iq_j \\pder{}{q_{i+j}} +\n\t\t\t2ij q_{i+j} \\pdertwo{}{p_i}{q_j} + ijp_{i+j} \\pdertwo{}{q_i}{q_j}\\\\\n\t\t\t+ \\frac{1}{2} (i+j) q_iq_j \\pder{}{p_{i+j}} + \\frac{1}{2} (i+j)\n\t\t\tp_ip_j \\pder{}{p_{i+j}} + ij p_{i+j} \\pdertwo{}{p_i}{p_j}\\biggr)\n\t\t\\end{multline}\n\t\tand\n\t\t\\begin{equation}\\label{Eq:CJ_2}\n\t\t\t\\mathcal{CJ}_2 = \\sum_{i=1}^\\infty \\left(i p_i\\frac{\\partial}{\\partial\n\t\t\t\tq_i} + i q_i\\frac{\\partial}{\\partial p_i}\\right)\n\t\t\\end{equation}\n\t\t\n\t\t\\subsection{Change of variables}\\label{SSec:Change}\n\t\t\n\t\tIn this section we reduce, by a suitable change of variables, the\n\t\toperators $\\mathcal{CJ}_1$ and $\\mathcal{CJ}_2$ to classical cut-and-join operators\n\t\t\\begin{equation}\\label{Eq:CJA}\n\t\t\t\\mathcal{CJ} \\stackrel{\\mbox{\\tiny def}}{=}\\sum_{i,j=1}^\\infty \\left(ijp_{i+j} \\pdertwo{}{p_i}{p_j}\n\t\t\t+(i+j) p_ip_j \\pder{}{p_{i+j}}\\right) \n\t\t\\end{equation}\n\t\tand Euler fields\n\t\t\\begin{equation}\\label{Eq:Euler}\n\t\t\tE\\stackrel{\\mbox{\\tiny def}}{=}\\sum_{i=1}^\\infty ip_i \\pder{}{p_i},\n\t\t\\end{equation}\n\t\trespectively.\n\t\t\n\t\t\\begin{proposition}\\label{Pp:ChVar}\n\t\t\tLet $u_\\ell = \\frac{p_\\ell+q_\\ell}{2}$ and $v_\\ell =\n\t\t\t\\frac{p_\\ell-q_\\ell}{2}$. Then\n\t\t\t\n\t\t\t\\begin{align*}\n\t\t\t\t\\mathcal{CJ}_1 &= \\sum_{i,j=1}^\\infty \\biggl( iju_{i+j}\\pdertwo{}{u_i}{u_j} +\n\t\t\t\t(i+j) u_iu_j \\pder{}{u_{i+j}} + ij v_{i+j} \\pdertwo{}{v_i}{v_j}\\\\\n\t\t\t\t&\\hphantom{\\sum_{i,j=1}^\\infty \\biggl(} + (i+j) v_iv_j \\pder{}{v_{i+j}}\\biggr) = \\mathcal{CJ}_u + \\mathcal{CJ}_v,\\\\\n\t\t\t\t\\mathcal{CJ}_2 &= \\sum_{\\ell=1}^\\infty \\ell \\left(u_\\ell \\pder{}{u_\\ell} -\n\t\t\t\tv_\\ell \\pder{}{v_{\\ell}}\\right) = E_u - E_v.\n\t\t\t\\end{align*}\n\t\t\\end{proposition}\n\t\t\\noindent where by $\\mathcal{CJ}_u$ and $\\mathcal{CJ}_v$ we denote the operator \\eqref{Eq:CJA}\n\t\twith $u_i$ (resp. $v_i$) substituted for $p_i$, and similarly, $E_u$\n\t\tand $E_v$.\n\t\t\n\t\tThe proof is a direct computation.\n\t\t\n\t\t\\begin{theorem}\\label{Th:B-Schur}\n\t\t\tFor all partitions $\\lambda$ and $\\mu$ the polynomials\n\t\t\t\\begin{equation*}\n\t\t\t\ts_{\\lambda\\mid\\mu}(p,q) \\stackrel{\\mbox{\\tiny def}}{=} s_\\lambda((p+q)\/2) s_\\mu((p-q)\/2)\n\t\t\t\\end{equation*} \n\t\t\twhere $s_\\lambda$ and $s_\\mu$ are Schur polynomials and $(p \\pm\n\t\t\tq)\/2$ means $((p_1 \\pm q_1)\/2, (p_2 \\pm q_2)\/2, \\dots)$, are\n\t\t\teigenvectors of both $\\mathcal{CJ}_1$ and $\\mathcal{CJ}_2$; the respective eigenvalues\n\t\t\tare $\\sum\\limits_{i=1}^\\infty \\left(\\lambda_i(\\lambda_i-2i+1) +\n\t\t\t\\mu_i(\\mu_i-2i+1)\\right)$ for $\\mathcal{CJ}_1$ and $\\sum\\limits_{i=1}^\\infty\n\t\t\t(\\lambda_i - \\mu_i)$ for $\\mathcal{CJ}_2$.\n\t\t\\end{theorem}\n\t\t\n\t\tThe theorem follows immediately from Proposition \\ref{Pp:ChVar} and\n\t\tthe classical fact that Schur polynomials are eigenvectors of the\n\t\tcut-and-join operator \\eqref{Eq:CJA} (see \\cite{LandoKazarian}) and of\n\t\tthe Euler field \\eqref{Eq:Euler} ($s_{\\lambda\\mid\\mu}$ are weighted\n\t\thomogeneous).\n\t\t\n\t\tTheorem \\ref{Th:B-Schur} and the Cauchy identity \\cite{Macdonald}\n\t\timply that\n\t\t\\begin{equation*}\n\t\t\te^{p_1} = e^{u_1} e^{v_1} = \\sum_{\\lambda,\\mu}\n\t\t\ts_{\\lambda\\mid\\mu}(p,q) s_{\\lambda\\mid\\mu}(1, 0, \\dots; 1, 0, \\dots)\n\t\t\\end{equation*}\n\t\t\n\t\tThis allows us to prove the following\n\t\t\n\t\t\\begin{corollary}[of Theorem \\ref{Th:B-Schur}]\\label{Cr:LKGenFunc}\n\t\t\t\\begin{align*}\n\t\t\t\t\\mathcal{H}^B(\\beta,\\gamma,p,q) &= \\sum_{\\lambda,\\mu}\n\t\t\t\t\\exp\\bigl(\\beta \\sum_{i=1}^\\infty (\\lambda_i(\\lambda_i-2i+1) +\n\t\t\t\t\\mu_i(\\mu_i-2i+1) + \\gamma \\sum_{i=1}^{\\infty}\n\t\t\t\t(\\lambda_i-\\mu_i)\\bigr)\\\\\n\t\t\t\t&\\times s_{\\lambda\\mid\\mu}(1,0,\\dots; 1,0,\\dots)\n\t\t\t\ts_{\\lambda\\mid\\mu}(p,q)\n\t\t\t\\end{align*}\n\t\t\\end{corollary}\n\t\tThis is the B-analog of the formula expressing the classical Hurwitz\n\t\tnumbers via the Schur polynomials \\cite{LandoKazarian}\n\t\t\n\t\tTheorem \\ref{Th:B-Schur} allows also to establish a direct relation\n\t\tbetween Hurwitz numbers of the groups $B_n$\n\t\t($h_{m,\\ell,\\lambda,\\mu}^B$, studied here) and $S_n$ (the classical\n\t\tones, denoted $h_{m,\\lambda}^A$ till the end of this section).\n\t\t\n\t\tNamely, it follows from equation \\ref{Eq:CJ_2} that $\\mathcal{CJ}_2(p_\\lambda\n\t\tq_\\mu)= (|\\lambda|+|\\mu|)p_\\lambda q_\\mu$, and therefore, by the\n\t\tcut-and-join equation,\n\t\t\\begin{multline*}\n\t\t\t\\sum_{m,\\ell}\\sum_{\\lambda,\\mu}\n\t\t\t\\frac{h_{m,\\ell,\\lambda,\\mu}^B}{m!(\\ell-1)!}\n\t\t\tp_{\\lambda}q_{\\mu}\\beta^m\\gamma^{\\ell-1} = \\mathcal{CJ}_2(\\mathcal{H}^B) =\n\t\t\t\\pder{\\mathcal{H}^B}{\\gamma}\\\\\n\t\t\t= \\sum_{m,\\ell}\\sum_{\\lambda,\\mu} (\\lmod\\lambda\\rmod +\n\t\t\t\\lmod\\mu\\rmod) \\frac{h_{m,\\ell,\\lambda,\\mu}^B}{m!\\ell!}\n\t\t\tp_{\\lambda}q_{\\mu}\\beta^m\\gamma^\\ell.\n\t\t\\end{multline*}\n\t\tThis implies the equality \n\t\t\\begin{equation*}\n\t\t\th_{m,\\ell,\\lambda,\\mu}^B = (\\lmod\\lambda\\rmod + \\lmod\\mu\\rmod)^\\ell\n\t\t\th_{m,0,\\lambda,\\mu}^B\n\t\t\\end{equation*}\n\t\tIt is thus enough to focus on $h_{m,0,\\lambda,\\mu}^B$.\n\t\t\n\t\tBy Corollary \\ref{Cr:LKGenFunc}, $h_{m,0,\\lambda,\\mu}^B$ is the\n\t\tcoefficient at the monomial $p_\\lambda q_\\mu\\beta^m$ in the expression\n\t\t\\begin{equation*}\n\t\t\t\\sum_{\\lambda',\\mu',m_1,m_2}\n\t\t\th_{m_1,\\lambda'}^A(p+q)_{\\lambda'} \n\t\t\th_{m_1,\\mu'}^A(p+q)_{\\mu'}\\beta^{m_1+m_2}\n\t\t\t\\frac{m!(\\lmod\\lambda\\rmod + \\lmod\\mu\\rmod)!}{2^{\\#\\lambda\n\t\t\t\t\t+\\#\\mu}\\lmod\\lambda'\\rmod! m_1! \\lmod\\mu'\\rmod! m_2!}\n\t\t\\end{equation*}\n\t\twhere $(p\\pm q)_\\lambda$ is understood as $(p_{\\lambda_1}\\pm\n\t\tq_{\\lambda_1}) \\dots (p_{\\lambda_s}\\pm q_{\\lambda_s})$. The\n\t\tcoefficient may be nonzero only if two conditions are satisfied:\n\t\t$m_1+m_2=m$ and $\\lambda'+\\mu'=\\lambda+\\mu$. Thus the relation can be\n\t\trewritten as\n\t\t\\begin{equation*}\n\t\t\th_{m,0,\\lambda,\\mu}^B =\n\t\t\t\\sum_{\\substack{\\lambda'+\\mu'=\\lambda+\\mu\\\\m_1+m_2=m}}\n\t\t\t\\frac{h_{m_1,\\lambda'}^Ah_{m_2,\\mu'}^A}{2^{\\#\\lambda +\\#\\mu}}\n\t\t\t\\binom{m}{m_1}\\binom{|\\lambda|+|\\mu|}{|\\lambda'|} [p_\\lambda q_\\mu\n\t\t\t:(p+q)_{\\lambda'}(p-q)_{\\mu'}]\n\t\t\\end{equation*}\n\t\t\n\t\tNow let $\\lambda'=(1^{\\alpha_1},2^{\\alpha_2},\\dots)$, $\\mu'=(1^{\\beta_1},2^{\\beta_2},\\dots)$ and\n\t\t$\\mu=(1^{\\gamma_1},2^{\\gamma_2},\\dots)$; it follows from the equality\n\t\t$\\lambda+\\mu = \\lambda' + \\mu'$ that $\\lambda =\n\t\t(1^{\\alpha_1+\\beta_1-\\gamma_1} 2^{\\alpha_2+\\beta_2-\\gamma_2} \\dots)$\n\t\tis fixed. If $f^{\\gamma}_{\\alpha\\beta}$ is a coefficient at\n\t\t$x^\\gamma$ at the polynomial $(1+x)^\\alpha(1-x)^\\beta$ then one has \n\t\t\n\t\t\\begin{equation*}\n\t\t\th_{m,0,\\lambda,\\mu}^B =\n\t\t\t\\sum_{\\substack{\\lambda'+\\mu'=\\lambda+\\mu\\\\m_1+m_2=m}}\n\t\t\t\\frac{h_{m_1,\\lambda'}^A h_{m_2,\\mu'}^A}{2^{\\#\\lambda + \\#\\mu}}\n\t\t\t\\binom{m}{m_1}\\binom{\\lmod\\lambda\\rmod +\n\t\t\t\t\\lmod\\mu\\rmod}{\\lmod\\lambda'\\rmod} f^{\\gamma_1}_{\\alpha_1\\beta_1}\n\t\t\tf^{\\gamma_2}_{\\alpha_2\\beta_2}\\dots\n\t\t\\end{equation*}\n\t\t\n\t\t\\section{Cut-and-join operator for the group $D_n$}\\label{Sec:CJD}\n\t\t\n\t\tLet $Z[D_n]$ be a center of the group algebra of $D_n$; it is a vector\n\t\tspace spanned by the class sums $\\CSum{\\lambda}{\\emptyset}^+$ and\n\t\t$\\CSum{\\lambda}{\\emptyset}^-$ where $\\lmod\\lambda\\rmod = n$ and all\n\t\tthe parts of $\\lambda$ are even, and the class sums\n\t\t$\\CSum{\\lambda}{\\mu}$ for all other $\\lambda$ and $\\mu$ such that \n\t\t$\\lmod\\lambda\\rmod + \\lmod\\mu\\rmod = n$ and $\\#\\mu$ is\n\t\teven. The operator of multiplication by the element $\\mathcal T_1 \\in\n\t\tZ[D_1]$ defined by \\eqref{Eq:DefT1} (a sum of all the reflections in\n\t\t$D_n$) acts on $Z[D_n]$; denote it $\\Theta_D$ for brevity.\n\t\t\n\t\tTake an element $x \\in B_n \\setminus D_n$ (i.e.\\ an odd permutation in\n\t\t$B_n$) and consider the map $\\mathcal{I}:\\mathbb{C}[D_n]\\mapsto\n\t\t\\mathbb{C}[D_n]$ defined by $\\mathcal{I}(a) = xax^{-1}$. Corollary\n\t\t\\ref{Cr:ConjOdd} implies that $\\mathcal I$ maps the center $Z[D_n]$ to\n\t\titself and its restriction to $Z[D_n]$ is an involution. This\n\t\tinvolution defines a splitting $Z[D_n] = V_n^+ \\oplus V_n^-$ where\n\t\t$V_n^{\\pm}$ are eigenspaces of $\\mathcal I$ corresponding to\n\t\teigenvalues $1$ and $-1$.\n\t\t\n\t\tClearly, the vector space $V_n^+$ is spanned by the elements\n\t\t$\\CSum{\\lambda}{\\mu}$ and $\\CSum{\\lambda}{\\emptyset}^++\n\t\t\\CSum{\\lambda}{\\emptyset}^-$, whereas the vector space $V_n^-$ is\n\t\tspanned by the the elements $\\CSum{\\lambda}{\\emptyset}^+-\n\t\t\\CSum{\\lambda}{\\emptyset}^-$. In particular, $V_n^-$ is nonempty if\n\t\tand only if $n = \\lmod\\lambda\\rmod$ is even. Therefore, there exist a\n\t\tlinear injection $\\Phi_+: V_n^+ \\to Z[B_n]$ and a linear\n\t\tisomorphism $\\Phi_-: V_n^- \\to Z[S_{n\/2}]$ (for $n$ even) defined as\n\t\t\\begin{align*}\n\t\t\t\\Phi_+(\\CSum{\\lambda}{\\mu}) &= \\CSum{\\lambda}{\\mu}\\\\\n\t\t\t\\Phi_+(\\CSum{\\lambda}{\\emptyset}^+ + \\CSum{\\lambda}{\\emptyset}^-) &= \\CSum{\\lambda}{\\emptyset},\n\t\t\\end{align*}\n\t\tand\n\t\t\\begin{equation*}\n\t\t\t\\Phi_-(\\CSum{\\lambda}{\\emptyset}^+ - \\CSum{\\lambda}{\\emptyset}^-) = C_{\\lambda\/2}\n\t\t\\end{equation*}\n\t\twhere $\\lambda\/2 = (\\lambda_1\/2 \\DT, \\lambda_s\/2)$ (recall that\n\t\t$\\CSum{\\lambda}{\\emptyset}^\\pm$ are defined only if all the parts\n\t\t$\\lambda_i$ of $\\lambda$ are even). The image of $\\Phi_+$ is a\n\t\tsubspace $Z_e[B_n] \\subset Z[B_n]$ spanned by all the class sums\n\t\t$\\CSum{\\lambda}{\\mu}$ with $\\#\\mu$ even.\n\t\t\n\t\t\\begin{theorem}\\label{Th:CJD}\n\t\t\t\\begin{enumerate}\n\t\t\t\t\\item\\label{It:Commutes} The operator $\\Theta_D$ commutes with the\n\t\t\t\tinvolution $\\mathcal I$, and therefore, $V_n^+$ and $V_n^-$ are\n\t\t\t\t$\\Theta_D$-invariant.\n\t\t\t\t\n\t\t\t\t\\item\\label{It:CJDonVPlus} For the restriction of $\\Theta_D$ on\n\t\t\t\t$V_n^+$ the following diagram is commutative:\n\t\t\t\t\\begin{equation}\\label{Eq:CJDViaGAlg}\n\t\t\t\t\t\\xymatrix{\n\t\t\t\t\t\tV_n^{+} \\ar[r]^{\\Theta_D} \\ar[d]^{\\Phi_+} & V_n^{+} \\ar[d]^{\\Phi_+} \\\\\n\t\t\t\t\t\tZ[B_n] \\ar[r]_{T_1} & Z[B_n]\n\t\t\t\t\t}\n\t\t\t\t\\end{equation}\n\t\t\t\t\n\t\t\t\t\\item\\label{It:CHDonVMinus} If $n$ is even then for the restriction\n\t\t\t\tof $\\Theta_D$ on $V_n^-$ the following diagram is commutative:\n\t\t\t\t\\begin{equation}\\label{Eq:CJDViaGAlgMinus}\n\t\t\t\t\t\\xymatrix{\n\t\t\t\t\t\tV_n^- \\ar[r]^{\\Theta_D} \\ar[d]^{\\Phi_-} & V_n^- \\ar[d]^{\\Phi_-} \\\\\n\t\t\t\t\t\tZ[S_{n\/2}] \\ar[r]_{4T} & Z[S_{n\/2}]\n\t\t\t\t\t}\n\t\t\t\t\\end{equation}\n\t\t\t\\end{enumerate}\n\t\t\\end{theorem}\n\t\t\n\t\t\\begin{corollary}\n\t\t\t$V_n^+$ is isomorphic to the space of homogeneous polynomials\n\t\t\t$\\mathbb{C}[p,q]_n^e$ of total degree $n$ and even degree in\n\t\t\t$q$; the isomorphism is given by \\eqref{Eq:DefIsoB}. $V_n^-$ is isomorphic to the space $\\mathbb{C}[p]_{n\/2}$ of\n\t\t\thomogeneous polynomials of total degree $n\/2$, the isomorphism is\n\t\t\tgiven by\n\t\t\t\\begin{equation*}\n\t\t\t\t\\Phi_A(\\CSum{\\lambda}{\\emptyset}^+-\\CSum{\\lambda}{\\emptyset}^-) =\n\t\t\t\tp_{\\lambda\/2} \\stackrel{\\mbox{\\tiny def}}{=} p_{\\lambda_1\/2} \\dots p_{\\lambda_s\/2}.\n\t\t\t\\end{equation*}\n\t\t\tThe following diagrams are commutative\n\t\t\t\\begin{equation*}\n\t\t\t\t\\xymatrix{\n\t\t\t\t\tV_n^+ \\ar[r]^{\\Theta_D} \\ar[d]^{\\Psi \\circ \\Phi_+} & V_n^+\n\t\t\t\t\t\\ar[d]^{\\Psi \\circ \\Phi_+} \\\\ \n\t\t\t\t\t\\mathbb{C}[p,q]_n^e \\ar[r]_{\\mathcal{CJ}_1} & \\mathbb{C}[p,q]_n^e\n\t\t\t\t}, \\qquad\n\t\t\t\t\\xymatrix{\n\t\t\t\t\tV_n^- \\ar[r]^{\\Theta_D} \\ar[d]^{\\Psi_A \\circ \\Phi_-} & V_n^-\n\t\t\t\t\t\\ar[d]^{\\Psi_A \\circ \\Phi_-} \\\\ \n\t\t\t\t\t\\mathbb{C}[p]_{n\/2} \\ar[r]_{\\mathcal{CJ}_A} & \\mathbb{C}[p]_{n\/2}\n\t\t\t\t}\n\t\t\t\\end{equation*}\n\t\t\twhere $\\mathcal{CJ}_A$ is the classical cut-and-join operator expressed in\n\t\t\trescaled variables $r_i = p_i\/2$ and multiplied by $4$. \n\t\t\\end{corollary}\n\t\t\n\t\tIn other words, the cut-and-join operator for the group $D_n$ is a\n\t\tdirect sum of the cut-and-join operator for the group $B_n$,\n\t\trestricted to the space of polynomials of even degree in $q$, and, for\n\t\t$n$ even, the cut-and-join operator for the group $S_{n\/2}$, expressed\n\t\tvia rescaled variables $r_i = p_i\/2$ and multiplied by $4$.\n\t\t\n\t\t\\begin{proof}[Proof of Theorem \\ref{Th:CJD} and the corollary]\n\t\t\tThe image of $V_n^+ \\subset Z[D_n] \\subset \\mathbb{C}[D_n]$ under the\n\t\t\tnatural embedding $\\iota: \\mathbb{C}[D_n] \\to \\mathbb{C}[B_n]$ is\n\t\t\t$Z[B_n]$, and the restriction of $\\iota$ to $V_n^+$ coincides with\n\t\t\t$\\Phi_+$. So, assertion \\ref{It:Commutes} follows from the fact that\n\t\t\t$\\Theta_D$ is an operator of multiplication by the element $T_1 \\in\n\t\t\tZ[D_n]$ such that $\\iota(T_1) \\in Z[B_n]$. Assertion\n\t\t\t\\ref{It:CJDonVPlus} is then evident.\n\t\t\t\n\t\t\tInstead of proving assertion \\ref{It:CHDonVMinus} we prove the\n\t\t\tcommutativity of the second diagram from the corollary; it is\n\t\t\tobviously an equivalent fact. The description above implies that\n\t\t\t$\\mathcal{CJ}_A = U \\mathcal{CJ} V$ where the operator $V: \\mathbb{C}[r] \\to\n\t\t\t\\mathbb{C}[p,q]$ is given by $(Vf)(p,q) = f(p_2, p_4, \\dots)$ and $U:\n\t\t\t\\mathbb{C}[p,q] \\to \\mathbb{C}[p]$, by $(Uh)(r) =\n\t\t\t\\left.h\\right|_{\\substack{p_1=p_3 \\DT= 0\\\\ p_2=r_1, p_4=r_2, \\dots\n\t\t\t\t\t\\\\ q_1=q_2 \\DT= 0}}$. Now the required formula follows directly\n\t\t\tfrom \\eqref{Eq:CJ_1} and \\eqref{Eq:CJA}.\n\t\t\\end{proof}\n\t\t\n\t\tHurwitz numbers for the group $D_n$ coincide with (some of) the\n\t\tHurwitz numbers for the group $B_n$ by Definition \\ref{Df:D-Hurwitz};\n\t\tsee also remark \\ref{Rm:DHurwRefine}.\n\t\t\n\t\t\\section{Other properties of B-Hurwitz numbers and B-cut-and-join operator}\n\t\t\n\t\tClassical Hurwitz numbers $h_{m,\\lambda}$ exhibit many interesting\n\t\tproperties. They play an active part in the algebro-geometric theory\n\t\tof moduli spaces of complex curves and holomorphic functions on them\n\t\t(see \\cite{GouldJacks,ELSV}), appear in some problems of\n\t\tlow-dimensional topology \\cite{BurmanZvonkine,TwistedHurwitz} and the\n\t\ttheory of integrable systems \\cite{Kramer,KazarianHodge,FSpace}. The\n\t\tclassical cut-and-join operator also has various interpretations\n\t\tbesides the one described here, including the Fock space\n\t\trepresentation (see \\cite{Johnson} and the references therein).\n\t\t\n\t\tIn this section we describe analogs of some of these results for the\n\t\tB-Hurwitz numbers and B-version of the cut-and-join operator. Proper\n\t\treferences and short explanation of the classical picture is given in\n\t\tthe beginning of each subsection.\n\t\t\n\t\t\\subsection{Ribbon decomposition}\\label{SSec:Ribbon}\n\t\t\n\t\tIn this section we will be using notations and definitions from\n\t\t\\cite{TwistedHurwitz}; to make the article self-contained, we quote\n\t\there the most important definitions from there. A {\\em\n\t\t\tdecorated-boundary surface} (DBS) with $n$ marked points is a\n\t\tsurface with the boundary $\\partial M$ and the marked points $a_1 \\DT, a_n\n\t\t\\in \\partial M$ satisfying some nontriviality conditions (see\n\t\t\\cite{TwistedHurwitz} for details). All DBS in this section are\n\t\tassumed oriented.\n\t\t\n\t\tIf $M$ is an oriented DBS with $n$ marked points then one defines its\n\t\t{\\em boundary permutation} $\\Sigma(M) \\in S_n$ as follows:\n\t\t$\\Sigma(M)(i) = j$ if the marked points $a_i$ and $a_j$ belong to the\n\t\tsame connected component of the boundary $\\partial M$ and $a_j$ immediately\n\t\tfollows $a_i$ as one moves around the boundary in the positive\n\t\tdirection (according to the orientation).\n\t\t\n\t\tA {\\em ribbon} is a long narrow rectangle $\\rho$ with alternating\n\t\tblack and white vertices; black vertices are joined by a diagonal. For\n\t\ta DBS $M$ and $1 \\le i < j \\le n$, denote by $G[i,j] M$ the result of\n\t\tgluing of a ribbon to $M$ such that the black vertices are identified\n\t\twith the marked points $a_i$ and $a_j$, and the short sides are glued\n\t\tto short segments of $\\partial M$ containing $a_i$ and $a_j$ and directed\n\t\taccording to the orientation of $\\partial M$. Apparently, $G[i,j] M$ is an\n\t\toriented DBS with the same marked points; we call $G[i,j]$ an\n\t\toperation of {\\em ribbon gluing}.\n\t\t\n\t\tA {\\em ribbon decomposition} of the DBS $M$ is an\n\t\torientation-preserving diffeomorphism of $M$ with the DBS $G[i_m,j_m]\n\t\t\\dots G[i_1,j_1] E_n$ obtained by consecutive ribbon gluing to a union\n\t\t$E_n$ of $n$ disks each one having one marked point $a_i$ on its\n\t\tboundary. The image in $M$ of the $k$-th ribbon glued is also called a\n\t\tribbon (number $k$) and is denoted $\\rho_k \\subset M$.\n\t\t\n\t\tNow fix positive integers $m$ and $\\ell$. Let $M$ be an oriented\n\t\tdecorated-boundary surface with $2n$ marked points $a_1 \\DT, a_{2n}$,\n\t\tand let $T: M \\to M$ be an orientation-preserving smooth involution\n\t\tmapping marked points to marked points and having no fixed points on\n\t\tthe boundary $\\partial M$. If $T(D) = D$ where $D \\subset \\partial M$ is a connected\n\t\tcomponent then $D$ contains an even number $2\\mu_i$ of marked\n\t\tpoints. If $D$ is not $T$-invariant then $T$ exchanges it with another\n\t\tcomponent $D'$; both contain the same number $\\lambda_i$ of marked\n\t\tpoints. Numbers $\\lambda_1 \\DT, \\lambda_s$ and $\\mu_1 \\DT, \\mu_t$ form\n\t\ttwo partitions denoted $\\Lambda_+(M,T)$ and $\\Lambda_-(M,T)$,\n\t\trespectively; one has $\\lmod\\Lambda_+(M,T)\\rmod +\n\t\t\\lmod\\Lambda_-(M,T)\\rmod = n$.\n\t\t\n\t\tLet now $R$ a ribbon decomposition of $M$ and suppose that it is\n\t\t$T$-invariant: for any ribbon $\\rho$ its image $T(\\rho)$ is also a\n\t\tribbon. Some ribbons are $T$-invariant: $T(\\rho) = \\rho$, while others\n\t\tcome in pairs $\\rho_1 = T(\\rho_2), \\rho_2 = T(\\rho_1)$ (for brevity,\n\t\tcall these $T$-pairs). If the diagonal of a ribbon $\\rho$ joins marked\n\t\tpoints $a_i$ and $a_j$ then the diagonal to the ribbon $T(\\rho)$ joins\n\t\t$T(a_i)$ and $T(a_j)$. In particular, diagonals of $T$-invariant\n\t\tribbons join marked points $a_i$ and $T(a_i)$, for some $i$.\n\t\t\n\t\tA ribbon is homeomorphic to a disk, so, by the Brouwer theorem, $T$\n\t\thas a fixed point on each $T$-invariant ribbon. Ribbons entering a\n\t\t$T$-pair do not contain fixed points: such ribbons, if not disjoint,\n\t\tintersect only at the boundary of $M$ where $T$ has no fixed points.\n\t\t\n\t\tFor any $i$, $1 \\le i \\le 2n$, denote by $\\bar i \\stackrel{\\mbox{\\tiny def}}{=} i + n \\bmod\n\t\t2n$. \n\t\t\n\t\t\\begin{definition}\\label{Def:B-RibDecom}\n\t\t\tA {\\em B-ribbon decomposition} with the profile\n\t\t\t$(\\lambda,\\mu,m,\\ell)$ is a DBS $M$ with $2n$ marked points equipped\n\t\t\twith a smooth orientation-preserving involution $T: M \\to M$ and a\n\t\t\t$T$-invariant ribbon decomposition $R$ such that\n\t\t\t\\begin{enumerate}\n\t\t\t\t\\item $T$ sends the marked point $a_i$ to the marked point $a_{\\bar\n\t\t\t\t\ti}$, for all $i = 1 \\DT, 2n$.\n\t\t\t\t\n\t\t\t\t\\item $\\Lambda_+(M,T) = \\lambda$, $\\Lambda_-(M,T) = \\mu$.\n\t\t\t\t\n\t\t\t\t\\item The number of $T$-invariant ribbons in the decomposition $R$\n\t\t\t\tis $\\ell$ and the number of $T$-pairs of ribbons is $m$.\n\t\t\t\t\n\t\t\t\t\\item $T$ sends the ribbon $\\rho_k$ either to itself or to the\n\t\t\t\tribbon $\\rho_{k\\pm1}$, for all $k = 1 \\DT, 2m+\\ell$.\n\t\t\t\t\n\t\t\t\t\\item $T$ has exactly one fixed point inside each $T$-invariant\n\t\t\t\tribbon and no fixed points elsewhere (in particular, it has no\n\t\t\t\tfixed points on the boundary of $M$).\n\t\t\t\\end{enumerate}\n\t\t\\end{definition}\n\t\t\n\t\tThe total space $M$ of a B-ribbon decomposition is thus a DBS with a\n\t\tribbon decomposition $H_1 H_2 \\dots H_{m+\\ell} E_{2n}$ where each\n\t\t$H_k$ is either an operation $G[i_k,\\bar i_k]$ of gluing a\n\t\t$T$-invariant ribbon or an operation $G[i,j] G[\\bar i,\\bar j]$ of\n\t\tgluing a $T$-pair.\n\t\t\n\t\tB-ribbon decompositions are split into equivalence classes in a\n\t\tnatural way: $(M_1,T_1,R_1)$ is said to be equivalent to\n\t\t$(M_2,T_2,R_2)$ if there exists an orientation-preserving\n\t\tdiffeomorphism $f: M_1 \\to M_2$ mapping marked points of $M_1$ to the\n\t\tmarked points of $M_2$ with the same numbers, each ribbon of $R_1$, to\n\t\tthe ribbon of $R_2$ with the same number, and transforming the\n\t\tinvolutions one into the other: $f \\circ T_1 = T_2 \\circ f$.\n\t\t\n\t\tLet $\\xi$ be a quotient of a ribbon $\\rho$ by the symmetry with\n\t\trespect to its center. We call $\\xi$ a {\\em petal}; it is a bigon with\n\t\ta black and a white vertex; one of its sides is long, and the other,\n\t\tshort. The image of the diagonal is a line joining the black vertex\n\t\twith an internal point of the petal.\n\t\t\n\t\tLet $M$ be a DBS with $n$ marked points, and $1 \\le i \\le n$. Denote\n\t\tby $G[i] M$ the result of gluing of a petal to $M$ such that the black\n\t\tvertex is identified with the marked point $a_i$, and the short side,\n\t\twith a short segment of $\\partial M$ containing $a_i$ and directed according\n\t\tto the orientation of $\\partial M$. Apparently, $G[i] M$ is a DBS, and its\n\t\tboundary inherits an orientation from $\\partial M$. We call $G[i]$ an\n\t\toperation of petal gluing; it is similar to the operation $G[i,j]$ of\n\t\tribbon gluing described above. \n\t\t\n\t\tLet $(M,T,R)$ be a B-ribbon decomposition with the profile\n\t\t$(\\lambda,\\mu,m,\\ell)$. Denote by $\\overline M$ the quotient of $M$ by\n\t\tthe involution $T$. Since $T$ has no fixed points on the boundary and\n\t\t$T(a_i) = a_{\\bar i}$, the quotient $\\overline M$ has a\n\t\tnatural structure of a DBS with $n$ marked points $b_i = p(a_i) =\n\t\tp(a_{\\bar i})$, $i = 1 \\DT, n$, where $p: M \\to \\overline M$ is\n\t\tthe natural projection. If the ribbons $\\rho_k$ and $\\rho_{k+1}$ form a\n\t\t$T$-pair then $\\rho = p(\\rho_1) = p(\\rho_2) \\subset \\overline M$ is a\n\t\tribbon; if the ribbon $\\rho_k$ is $T$-invariant then $p(\\rho_k)\n\t\t\\subset \\overline M$ is a petal. In other words, if $M = H_1 \\dots\n\t\tH_{m+\\ell} E_{2n}$ is a ribbon decomposition described above then\n\t\t$\\overline M = \\overline H_1 \\dots \\overline H_{m+\\ell} E_n$ where\n\t\t$\\overline H_k = G[i_k,j_k]$ if $H_k = G[i_k, j_k] G[\\bar i_k, \\bar\n\t\tj_k]$ and $\\overline H_k = G[i_k]$ if $H_k = G[i_k, \\bar i_k]$. \n\t\t\n\t\t\\begin{theorem}\\label{Th:1-1Corresp}\n\t\t\tLet $\\lambda$ and $\\mu$ be two partitions such that\n\t\t\t$\\lmod\\lambda\\rmod + \\lmod\\mu\\rmod = n$, and $m$, $\\ell$, two\n\t\t\tpositive integers. There is a one-to-one correspondences between\n\t\t\tsequences of reflections in the group $B_n$ having profile\n\t\t\t$(\\lambda,\\mu,m,\\ell)$ and equivalence classes of B-ribbon\n\t\t\tdecompositions having the same profile.\n\t\t\\end{theorem}\n\t\t\n\t\t\\begin{proof}\n\t\t\tThe correspondence in question relates a sequence of reflections $h_1\n\t\t\t\\DT, h_{m+\\ell}$ to a ribbon decomposition $H_1 \\dots H_{m+\\ell}\n\t\t\tE_{2n}$ of the DBS $M$ where for every $k$ one has $H_k = G[i,j]G[\\bar\n\t\t\ti, \\bar j]$ if $h_k = (ij)(\\bar i \\bar j)$ and $H_k = G[i,\\bar i]$\n\t\t\tif $h_k = (i \\bar i)$. The involution $T$ exchanges the two ribbons\n\t\t\tattached in the first case and leaves the ribbon invariant in the\n\t\t\tsecond case. By \\cite[Theorem 1.11]{TwistedHurwitz} the boundary\n\t\t\tpermutation of $M$ is equal to $h_1 \\dots h_{m+\\ell} \\in B_n$. It\n\t\t\tfollows from the description of the conjugacy classes in Section\n\t\t\t\\ref{SSec:ConjClass} that $h_1 \\dots h_{m+\\ell} \\in\n\t\t\tC_{\\lambda\\mid\\mu}$ if and only if $\\Lambda_+(M,T) = \\lambda$ and\n\t\t\t$\\Lambda_-(M,T) = \\mu$.\n\t\t\\end{proof}\n\t\t\n\t\t\\subsection{Fermionic interpretation of B-cut-and-join}\\label{SSec:Fermion}\n\t\t\n\t\tFirst, recall some classical results about Fock space,\n\t\ta.k.a.\\ (semi-)infinite wedge product. For a partition $\\lambda =\n\t\t(\\lambda_1 \\DT\\ge \\lambda_s)$ define a {\\em Fock set}\n\t\t\\begin{equation*}\n\t\t\t\\Phi_\\lambda \\stackrel{\\mbox{\\tiny def}}{=} \\{i - \\frac12 - \\lambda_i \\mid i = 1, 2, \\dots\\}\n\t\t\t\\subset \\Integer + \\frac12;\n\t\t\\end{equation*}\n\t\there one assumes $\\lambda_i = 0$ if $i > s$.\n\t\t\n\t\t\\begin{lemma}\n\t\t\t\\begin{enumerate}\n\t\t\t\t\\item $\\#\\bigl(\\Phi_\\lambda \\cap \\{x \\mid x < 0\\}\\bigr) =\n\t\t\t\t\\#\\bigl(\\{x \\mid x > 0\\} \\setminus \\Phi_\\lambda =\n\t\t\t\t\\lmod\\lambda\\rmod \\stackrel{\\mbox{\\tiny def}}{=} \\lambda_1 \\DT+ \\lambda_s$.\n\t\t\t\t\n\t\t\t\t\\item If a subset $\\Phi \\subset \\Integer + \\frac12$ is such that $\\Phi \\cap \\{x\n\t\t\t\t\\mid x < 0\\}$ and $\\{x \\mid x > 0\\} \\setminus \\Phi$ are finite sets\n\t\t\t\tof equal cardinality then there exists a unique partition $\\lambda$\n\t\t\t\tsuch that $\\Phi = \\Phi_\\lambda$.\n\t\t\t\\end{enumerate}\n\t\t\\end{lemma}\n\t\t\n\t\tThe lemma is proved by direct induction; see \\cite{FSpace} for\n\t\tdetails.\n\t\t\n\t\tA classical (charge zero) Fock space \\cite{FSpace} (or fermion state\n\t\tspace) is a vector space $\\Lambda^{\\infty\/2}$ with the basis $w_\\lambda$ where\n\t\t$\\lambda$ runs through all partitions, including the empty partition;\n\t\t$w_\\emptyset$ is called the vacuum vector. One usually represents\n\t\t$w_\\lambda$ as an ``infinite wedge product'' $\\bigwedge_{u \\in\n\t\t\t\\Phi_\\lambda} z_u = z_{1\/2-\\lambda_1} \\wedge z_{3\/2-\\lambda_2}\n\t\t\\wedge \\dots$; note that the indices at the product form an increasing\n\t\tsequence.\n\t\t\n\t\tLet $k \\in \\mathbb{Z}$, $k \\ne 0$. Consider a linear operator $a_k:\n\t\t\\Lambda^{\\infty\/2} \\to \\Lambda^{\\infty\/2}$ acting on the basic vectors as follows:\n\t\t\\begin{equation*}\n\t\t\ta_k (\\bigwedge_{u \\in \\Phi_\\lambda} z_u) = \\sum_{v \\in \\Phi_\\lambda}\n\t\t\tz_{v-k} \\wedge (-1)^{d_{\\lambda,k}(v)} \\bigwedge_{u \\in \\Phi_\\lambda\n\t\t\t\t\\setminus \\{v\\}} z_u\n\t\t\\end{equation*}\n\t\twhere $d_{\\lambda,k}(v)$ is defined as $\\#\\{u \\in \\Phi_\\lambda \\mid\n\t\tv-k \\le u < v\\}$ if $k > 0$ and as $\\#\\{u \\in \\Phi_\\lambda \\mid v < u\n\t\t\\le v-k\\}$ if $k < 0$.\n\t\t\n\t\t\\begin{theorem}[\\cite{FSpace}]\\label{Th:Fock}\n\t\t\t\\begin{enumerate}\n\t\t\t\t\\item Operators $a_i$ and $a_j$ commute unless $i+j = 0$.\n\t\t\t\t\n\t\t\t\t\\item\\label{It:BFC} For any element $x \\in \\Lambda^{\\infty\/2}$ there exists a\n\t\t\t\tunique polynomial $R_x \\in \\mathbb{C}[p] \\stackrel{\\mbox{\\tiny def}}{=} \\mathbb{C}[p_1, p_2,\n\t\t\t\t\\dots]$ such that $x = R_x(a_1,a_2,\\dots)w_{\\emptyset}$. If $x =\n\t\t\t\tw_\\lambda$ then $R_x = s_\\lambda$ (a classical Schur polynomial).\n\t\t\t\\end{enumerate}\n\t\t\\end{theorem}\n\t\t\n\t\tThe map $R: \\Lambda^{\\infty\/2} \\to \\mathbb{C}[p]$ defined in \\ref{It:BFC} is called\n\t\tboson-fermion correspondence. It is a linear bijection, so for any\n\t\tlinear operator $A: \\Lambda^{\\infty\/2} \\to \\Lambda^{\\infty\/2}$ one can define its polynomial\n\t\tcounterpart $\\hat A \\stackrel{\\mbox{\\tiny def}}{=} R A R^{-1}: \\mathbb{C}[p] \\to \\mathbb{C}[p]$;\n\t\tthe map $A \\leftrightarrow \\hat A$ is an isomorphism of the algebras\n\t\tof linear operators on $\\Lambda^{\\infty\/2}$ and on $\\mathbb{C}[p]$.\n\t\t\n\t\t\\begin{proposition}[\\cite{FSpace}]\\label{Pp:OperHat}\n\t\t\t\\begin{enumerate}\n\t\t\t\t\\item\\label{It:PdP} If $A = a_k$ then $\\hat A = p_k$ (the operator of\n\t\t\t\tmultiplication) if $k > 0$ and $\\hat A = -k \\pder{}{p_{-k}}$ if $k\n\t\t\t\t< 0$.\n\t\t\t\t\n\t\t\t\t\\item\\label{It:Energy} If $A = \\sum_{i=1}^\\infty a_i a_{-i}$ (called\n\t\t\t\tthe energy operator) then $\\hat A$ is the Euler vector field\n\t\t\t\t$\\sum_{i=1}^\\infty i p_i \\pder{}{p_i}$.\n\t\t\t\t\n\t\t\t\t\\item\\label{It:CJ} If $A = \\sum_{\\substack{i,j,k \\ne 0\\\\ i+j+k = 0}}\n\t\t\t\ta_i a_j a_k$ then $\\hat A$ is the classical cut-and-join operator\n\t\t\t\t$\\mathcal{CJ}$.\n\t\t\t\\end{enumerate}\n\t\t\\end{proposition}\n\t\t\n\t\tIt can be proved (see \\cite{FSpace}) that for any partition $\\lambda$\n\t\tthe basic element $w_\\lambda$ is an eigenvector of both the energy\n\t\toperator and the operator from assertion \\ref{It:CJ} of Proposition\n\t\t\\ref{Pp:OperHat}. The corresponding eigenvalues are $\\lmod\n\t\t\\lambda\\rmod \\stackrel{\\mbox{\\tiny def}}{=} \\sum_i \\lambda_i$ and $\\phi(\\lambda) \\stackrel{\\mbox{\\tiny def}}{=} \\sum_i\n\t\t\\lambda_i(\\lambda_i-2i+1)$. Equivalently, the Schur polynomial\n\t\t$s_\\lambda$ is the eigenvector of the Euler field with the eigenvalue\n\t\t$\\lmod \\lambda\\rmod = \\deg s_\\lambda$ and of the cut-and-join operator\n\t\twith the eigenvalue $\\phi(\\lambda)$.\n\t\t\n\t\tConsider now a space $L \\stackrel{\\mbox{\\tiny def}}{=} \\Lambda^{\\infty\/2} \\otimes \\Lambda^{\\infty\/2}$ (two-particle\n\t\tfermion state space) with the basis $w_\\lambda \\otimes w_\\mu$, where\n\t\t$\\lambda, \\mu$ are partitions. The tensor square of the boson-fermion\n\t\tcorrespondence $R^{\\otimes 2}: L \\to \\mathbb{C}[p]^{\\otimes 2}\n\t\t=\\mathbb{C}[u,v]$ where $u = (u_1, u_2, \\dots)$ and $v = (v_1, v_2,\n\t\t\\dots)$ is a linear isomorphism. Like for the classical case, any\n\t\toperator $A: L \\to L$ has its polynomial counterpart $R^{\\otimes 2} A\n\t\t(R^{\\otimes 2})^{-1}: \\mathbb{C}[u,v] \\to \\mathbb{C}[u,v]$, which we also\n\t\tdenote by $\\hat A$.\n\t\t\n\t\tFor every $k \\ne 0$ consider two linear operators: $b_k \\stackrel{\\mbox{\\tiny def}}{=} a_k\n\t\t\\otimes I + I \\otimes a_k$ and $c_k \\stackrel{\\mbox{\\tiny def}}{=} a_k \\otimes I - I \\otimes\n\t\ta_k$. It follows from Proposition \\ref{Pp:OperHat} (assertion\n\t\t\\ref{It:PdP}) that\n\t\t\\begin{align}\n\t\t\t&\\begin{aligned}\n\t\t\t\t&\\hat b_k = u_k + v_k = p_k,\\\\\n\t\t\t\t&\\hat c_k = u_k - v_k = q_k\n\t\t\t\\end{aligned}\n\t\t\t\\qquad \\text{if } k > 0,\\label{Eq:KPos}\\\\\n\t\t\t&\\begin{aligned}\n\t\t\t\t&\\hat b_k = -k(\\pder{}{u_{-k}} + \\pder{}{v_{-k}}) =\n\t\t\t\t-2k\\pder{}{p_{-k}},\\\\\n\t\t\t\t&\\hat c_k = -k(\\pder{}{u_{-k}} - \\pder{}{v_{-k}}) =\n\t\t\t\t-2k\\pder{}{q_{-k}}\n\t\t\t\\end{aligned}\n\t\t\t\\qquad \\text{if } k < 0.\\label{Eq:KNeg}\n\t\t\\end{align}\n\t\tHere we are using notation from Section \\ref{SSec:Change}.\n\t\t\n\t\tApparently, $b_i$ and $b_j$, as well as $c_i$ and $c_j$, commute unless\n\t\t$i+j = 0$; $b_i$ and $c_j$ commute for all $i, j$.\n\t\t\n\t\t\\begin{proposition}\n\t\t\tIf \n\t\t\t\\begin{align*}\n\t\t\t\tJ_1 &\\stackrel{\\mbox{\\tiny def}}{=} \\frac{1}{12}(\\sum_{\\substack{i+j+k=0\\\\i,j,k \\ne 0}}\n\t\t\t\tb_ib_jb_k + \\sum_{\\substack{i+j+k=0\\\\i,j,k \\ne 0}}b_ic_jc_k +\n\t\t\t\t\\sum_{\\substack{i+j+k=0\\\\i,j,k\\neq0}}\n\t\t\t\tc_ib_jc_k+\\sum_{\\substack{i+j+k=0\\\\i,j,k \\ne 0}}c_ic_jb_k)\\\\\n\t\t\t\t&= \\frac{1}{12}\\sum_{\\substack{i+j+k=0\\\\i,j,k\\neq0}} b_ib_jb_k +\n\t\t\t\t\\frac{1}{4} \\sum_{\\substack{i+j+k=0\\\\i,j,k \\ne 0}} b_ic_jc_k,\\\\\n\t\t\t\tJ_2 &\\stackrel{\\mbox{\\tiny def}}{=} \\sum_{\\ell=1}^\\infty b_\\ell c_{-\\ell} +\n\t\t\t\t\\sum_{\\ell=1}^\\infty c_\\ell b_{-\\ell}\n\t\t\t\\end{align*}\n\t\t\tthen $\\hat{J_1} = \\mathcal{CJ}_1$ and $\\hat{J_2} = \\mathcal{CJ}_2$ (the B-cut-and-join\n\t\t\toperators), and also\n\t\t\t\\begin{align*}\n\t\t\t\tJ_1 &= C \\otimes \\name{Id} + \\name{Id} \\otimes C,\\\\\n\t\t\t\tJ_2 &= E\\otimes \\name{Id} - \\name{Id}\\otimes E\n\t\t\t\\end{align*}\n\t\t\twhere $C$ and $E$ are the operators defined in Proposition\n\t\t\t\\ref{Pp:OperHat}.\n\t\t\\end{proposition}\n\t\t\n\t\tThe proof is straightforward: the first pair of equations is proved by\n\t\ta calculation using \\eqref{Eq:KPos}, \\eqref{Eq:KNeg}, \\eqref{Eq:CJ_1}\n\t\tand \\eqref{Eq:CJ_2}, as well as the commutativity relations. The\n\t\tsecond pair of equations follows from Proposition \\ref{Pp:ChVar} or\n\t\tcan be obtained by a direct calculation, too.\n\t\t\n\t\t\\begin{corollary}\\label{Th:Fock2}\n\t\t\t$w_\\lambda\\otimes w_\\mu$ are eigenvectors of $J_1$ and $J_2$ with the\n\t\t\teigenvalues $\\sum_{i=1}^\\infty (\\lambda_i(\\lambda_i-2i+1) +\n\t\t\t\\mu_i(\\mu_i-2i+1)) $ and $\\sum_{i=1}^\\infty (\\lambda_i-\\mu_i)$,\n\t\t\trespectively.\n\t\t\\end{corollary}\n\t\t\n\t\t\\subsection{KP hierarchy for B-Hutwitz numbers}\n\t\t\n\t\tThe KP hierarchy is one of the best known and studied integrable\n\t\tsystems; see e.g.\\ \\cite{Kramer,KazarianHodge,FSpace} for details. It\n\t\tis an infinite system of PDE applied to a formal series $F \\in\n\t\t\\mathbb{C}[[t]]$ where $t = (t_1, t_2, \\dots)$ is a countable collection\n\t\tof variables (``times''). Each equation looks like\n\t\t$\\pdertwo{F}{t_i}{t_j} = P_{ij}(F)$ where $P_{ij}$ is a homogeneous\n\t\tmultivariate polynomial with partial derivatives of $F$ used as its\n\t\targuments. For example, the first two equations of the hierarchy are\n\t\t\n\t\t\\begin{align*}\n\t\t\t&F_{22}=-\\frac{1}{2}F_{11}^2 + F_{31} -\\frac{1}{12}F_{1111}\\\\\n\t\t\t&F_{32}=-F_{11}F_{21}+F_{41}-\\frac{1}{6}F_{2111};\n\t\t\\end{align*}\n\t\there $F_{i_1i_2\\dots i_n}$ means $\\frac{\\partial^nF}{\\partial\n\t\t\tu_{i_1}\\partial u_{i_2\\dots\\partial u_{i_n}}}$. If $F$ is a solution\n\t\tof the hierarchy, its exponential $\\tau= e^F$ is called a\n\t\t$\\tau$-function.\n\t\t\n\t\t\\begin{example}\\label{Ex:T1Tau}\n\t\t\tEquations of the hierarchy do not contain first-order derivatives,\n\t\t\tso the function $F(t) = t_1$ is a solution of the KP hierarchy;\n\t\t\t$\\tau = e^{t_1}$ is a $\\tau$-function.\n\t\t\\end{example}\n\t\t\n\t\t\\begin{proposition}\\label{Pp:Sato}\n\t\t\tThere exists a Lie algebra $\\mathcal G$ of differential operators on\n\t\t\tthe space $\\mathbb{C}[[p]]$ having the following properties:\n\t\t\t\\begin{enumerate}\n\t\t\t\t\\item\\label{It:GLInf} $\\mathcal G$ is isomorphic to a\n\t\t\t\tone-dimensional central extension $\\widehat{\\name{gl}(\\infty)}$ of\n\t\t\t\ta Lie algebra $\\name{gl}(\\infty)$ of infinite matrices\n\t\t\t\t$(x_{ij})_{i,j \\in \\mathbb{Z}}$ having nonzero elements at a finite\n\t\t\t\tset of diagonals.\n\t\t\t\t\n\t\t\t\t\\item\\label{It:F2F3} The Euler field $E$ and the cut-and-join\n\t\t\t\toperator $\\mathcal{CJ}$ belong to $\\mathcal G$.\n\t\t\t\t\n\t\t\t\t\\item\\label{It:Orbit} For any $\\tau$-function $f$ and any $X \\in\n\t\t\t\t\\mathcal G$, the power series $e^X f$ is a $\\tau$-function,\n\t\t\t\ttoo. \n\t\t\t\\end{enumerate}\n\t\t\\end{proposition}\n\t\t\n\t\tSee \\cite{Kramer} for the definition of $\\mathcal G$, \\cite{2Sato} and\n\t\t\\protect{\\cite[section 6]{KazarianHodge}} for the proof.\n\t\t\n\t\t\\begin{theorem}\\label{Th:HBsolKP}\n\t\t\tThe generating function $\\mathcal{H}^B(\\beta,\\gamma,u+v,u-v)$ is a\n\t\t\t$2$-parameter family of $\\tau$-functions, independently in the $u$\n\t\t\tand the $v$ variables.\n\t\t\\end{theorem} \n\t\t\n\t\t\\begin{proof}\n\t\t\tEquation \\eqref{Eq:GenFunc} and Proposition \\ref{Pp:ChVar} imply that\n\t\t\t\\begin{align*}\n\t\t\t\t\\mathcal{H}^B(\\beta,\\gamma,u+v,u-v) = e^{\\beta(\\mathcal{CJ}_u+\n\t\t\t\t\t\\mathcal{CJ}_v)}e^{\\gamma(E_u-E_v)} e^{u_1}e^{v_1} =\n\t\t\t\te^{(\\beta\\mathcal{CJ}_{u}+\\gamma E_u)}e^{u_1} e^{(\\beta\\mathcal{CJ}_{v}-\\gamma\n\t\t\t\t\tE_v)}e^{v_1}\n\t\t\t\\end{align*}\n\t\t\t(operators $\\mathcal{CJ}_i$ and $\\mathcal{CJ}_v$ commute because they act on different\n\t\t\tsets of variables).\n\t\t\t\n\t\t\tBy assertion \\ref{It:F2F3} of Proposition \\ref{Pp:Sato},\n\t\t\t$\\beta\\mathcal{CJ}_{u}+ \\gamma E_u \\in \\mathcal G_u$ (the Lie algebra $\\mathcal\n\t\t\tG$ acting by differential operators on the $u$ variables). By Example\n\t\t\t\\ref{Ex:T1Tau}, $e^{u_1}$ is a $\\tau$-function (in the $u$ variables),\n\t\t\tand so is $\\mathcal{H}^B(\\beta,\\gamma,u+v,u-v)$ by assertion\n\t\t\t\\ref{It:Orbit} of the same proposition. The reasoning for the $v$\n\t\t\tvariables is the same.\n\t\t\\end{proof}\n\t\t\n\t\t\\begin{corollary}\n\t\t\tThe generating function $\\mathcal{H}^D(\\beta,u+v,u-v)$ is a 1-parameter\n\t\t\tfamily of $\\tau$-functions independently in the $u$ and the $v$\n\t\t\tvariables.\n\t\t\\end{corollary}\n\t\t\\begin{proof}\n\t\t\tWe apply the same strategy as in the proof of theorem\n\t\t\t\\ref{Th:HBsolKP} to the function $\\mathcal{H}^D(\\beta,u+v,u-v) =\n\t\t\te^{\\beta(\\mathcal{CJ}_u+ \\mathcal{CJ}_v)}e^{u_1}e^{v_1}$.\n\t\t\\end{proof}\n\t\t\n\t\t","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{INTRODUCTION}\n\nHydrogen and helium are the two most abundant elements in giant\nplanets. While Jupiter and Saturn are well characterized on the\nsurface, many basic questions about its interior have not been\nanswered\\cite{Gu02}. Jupiter's surface composition has been measured\n{\\em in situ} by the Galileo probe:\\cite{FH03} H 74.2\\% by weight, He\n23.1\\%, and 0.027\\% heavier elements, which is enhanced compared to\nthe protosolar composition: 0.015\\% heavier elements along with H\n73.6\\% and He 24.9\\%. The abundance of heavy element in the interior\nand their distribution are not well characterized.\\cite{SG04} In\nparticular, it is conversial whether Jupiter has a rocky core. The\ndetection of a core in Jupiter may validate the standard model of\ngiant-planet formation, nucleated capture of nebular\nhydrogen.\\cite{Li93} An alternative scenario was proposed by Alan Boss\nwho suggested that giant planets form directly from spiral\ninstabilities in protostellar disks.\\cite{Boss97,Boss04} Under this\ngravitational instability hypothesis, giant planets would not have a\ncore, or at least a much smaller one.\n\nSince there is no direct way to detect a core in Jupiter, one must\ninstead refer to indirect measurements and to models for the planet\ninterior. Such models are constrained by the available observation\ndata, in particular, the properties at the planet surface and the\ngravitational moments measured through fly-by trajectories. All models\nrely on an equation of state (EOS) of hydrogen-helium\nmixtures. However, the uncertainties in the available EOS are large\nand have not allowed one, among many other questions, to determine\nwhether Jupiter has a core.\n\nIn this article, we present results from path integral Monte Carlo\n(PIMC) simulations\\cite{Ce95} that enable us to study quantum many-body\nsystems at finite temperature from first principles. In this\nsimulation, hydrogen-helium mixtures are represented by an ensemble of\nelectrons, protons and helium nuclei, each described by a path in\nimaginary time to incorporate quantum effects. The electrons are\ntreated as fermions while exchange effects for the nuclei can be\nneglected for the considered thermodynamics conditions.\n\n\\section{PATH INTEGRAL MONTE CARLO}\n\nThe thermodynamic properties of a many-body quantum system at finite\ntemperature can be computed by averaging over the density matrix,\n$\\hat{\\rho} = e^{-\\beta \\hat{H}}, \\beta=1\/k_{\\rm B} T$. \nPath integral formalism\\cite{Fe53} is based on the identity,\n\\begin{equation}\ne^{-\\beta \\hat{H}} = \\left[ e^{-\\frac{\\beta}{M} \\hat{H}} \\right]^M\n\\end{equation}\nwhere $M$ is a positive integer. Insertion of complete sets of states\nbetween the $M$ factors leads to the usual imaginary time path\nintegral formulation, written here in real space,\n\\begin{equation}\n\\rho({\\bf R},{\\bf R}';\\beta)=\n\\int\\ldots\\int d{\\bf R}_{1}\\ldots d{\\bf R}_{M-1} \\; \\rho({\\bf R},{\\bf R}_{1};\\tau)\\ldots\\rho({\\bf R}_{M-1},{\\bf R}';\\tau)\n \\label{eq2.5}\n\\end{equation}\nwhere $\\tau=\\beta\/M$ is the time step. Each of the $M$ steps in the\npath now has a high temperature density matrix\n$\\rho({\\bf R}_k,{\\bf R}_{k+1};\\tau)$ associated with it. The integrals are\nevaluated by Monte Carlo methods. \nThe density matrix for bosonic and fermionic systems can be obtained\nby projecting out states of corresponding symmetry. In PIMC, one sums\nup different permutations $P$,\n\\begin{equation}\n\\nonumber\n\\rho_{\\rm b\/f}({\\bf R},{\\bf R}';\\beta) \n= \n\\frac{1}{N!} \\sum_{\\mathcal{P}} \\; (\\pm 1)^{\\mathcal{P}} \\rho({\\bf R},{\\mathcal{P}} {\\bf R}' ; \\beta)\n= \\frac{1}{N!} \\sum_{\\mathcal{P}} \\; (\\pm 1)^{\\mathcal{P}} \\; \\! \\! \\! \\! \\! \\! \\! \\int \\limits_{{\\bf R} \\rightarrow {\\mathcal{P}} {\\bf R}' } \n\\! \\! \\! \\! \\! \\! d{\\bf R}_t \\, e^{-U[{\\bf R}_t] }.\n\\end{equation}\nFor bosons, this is essentially an exact numerical algorithm that has\nfound many applications.\\cite{Ce95} For fermions, the cancellation of\npositive and negative contributions leads to numerically unstable\nmethods, which is known as the {\\em fermion sign problem}. Ceperley\nshowed that this problem can be solved by introducing the {\\em\nrestricted path} approximation,\\cite{Ce91,Ce96}\n\\begin{equation} \n\\rho_f({\\bf R}_0, {\\bf R}' ;\\beta) \\approx\n\\frac{1}{N!}\\; \\sum_{\\mathcal{P}} \\; (-1)^{\\mathcal{P}} \n\\int\\limits_{\\scriptsize {\n\\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\!\\! \\! \\! \\! \n\\begin{array}{c} \n {\\bf R}_0 \\rightarrow {\\mathcal{P}} {\\bf R}'\\\\\n {\\rho_T({\\bf R}(t),{\\bf R}_0;t)>0}\n\\end{array}}\n \\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\! \\!\n\\! \\! \\! \\! \\! \\! \\! \\! \\! \\! } d{\\bf R}_t \\;\\; e^{-U[{\\bf R}(t)] } \\quad,\n\\label{restricted_PI} \n\\end{equation}\nwhere one only samples path ${\\bf R}(t)$ that stay within the positive\nregion of a {\\em trial} density matrix,\n$\\rho_T({\\bf R}(t),{\\bf R}_0;t)>0$. This procedure leads to an efficient\nalgorithm for fermionic systems. All negative contributions to\ndiagonal matrix elements are eliminated.\\cite{MPC03} Contrary to the\nbosonic case, the algorithm is no longer exact since it now depends on\nthe approximations for the trial density matrix. However, the method\nhas worked very well in many applications.\\cite{PC94,MP02} For\n$\\rho_T$, one can use the free particle or a variational density\nmatrix.\\cite{MP00}\nIn the results presented here, the\nhigh temperature density matrix was taken as a product of exact pair\ndensity matrices,\n\\begin{eqnarray}\n\\frac{\\rho({\\bf R},{\\bf R}';\\tau)}{\\rho_0({\\bf R},{\\bf R}';\\tau)} \n&=& \n\\left< e^{-\\int_0^\\tau dt \\sum_{i_{\\rm {\\bf R} \\to {\\bf R}'}\n=\n\\left< \\prod_{i_{\\rm {\\bf R} \\to {\\bf R}'}\\\\\n&\\approx&\n\\prod_{i_{\\rm {\\bf r}_{ij} \\to {\\bf r}_{ij}'} \n\\equiv\ne^{ - \\sum_{i$. This approach is\nefficient but not exact, and therefore puts a limit on the imaginary\ntime step $\\tau$ in many-body simulations. The pair action, $u$, can\nbe computed by different methods.\\cite{St68,Po88,KS95}\nThe following simulation results were derived with $N_e=32$ electrons\nand the corresponding number of protons, $N_p$, and helium nuclei\n$N_{\\rm He}$ to obtain a neutral system ($N_e=N_p+2 N_{\\rm He}$) with\nhelium fraction $x\\equiv 2 N_{\\rm He}\/N_e$. Periodic boundary conditions\nare applied. As imaginary time discretization, we employ\n$\\tau=0.079$. (Atomic units of Bohr radii and Hartrees will be used\nthroughout this article.)\n\nThe electrons are treated as fermions with fixed spin. We use\nvariational nodes to restrict the paths and have therefore extended\nthe approach in Ref.~[\\onlinecite{MP00}] to mixtures. The nuclei\nare also treated as paths but their exchange effects are not relevant\nhere.\n\n\\section{PHASE DIAGRAM OF HOT DENSE HYDROGEN}\n\n\\begin{figure}[htb]\n\\centerline{\\includegraphics[angle=0,width=4in]{log_astro_phase_diagram12.eps}}\n\\caption{Density-temperature phase diagram of hot dense hydrogen. \n The blue dash-dotted lines separate the molecular, the atomic,\n the metallic, and the plasma regime. The green solid lines are\n isentropes for Jupiter and stars with 0.3, 1, and 15 solar\n masses. Single shock Hugoniot states as well as the inertial\n confinement fusion path\\cite{lindl} are indicated by dashed\n lines. The thin solid line show $\\rho$-$T$ conditions of PIMC\n simulations.}\n\\label{phase}\n\\end{figure}\n\nFigure~\\ref{phase} shows the high temperature phase diagram of dense\nhydrogen beginning with the fluid and reaching up to a highly ionized\nplasma state. The figure includes the isentropes for Jupiter and low\nmass stars\\cite{SC95} and indicates the thermodynamic conditions, at\nwhich PIMC simulations have been applied.\\cite{MC00,MC01,MP02} We are\nnow going to use these simulation results to characterize hot dense\nhydrogen from a path integral perspective.\n\n\\begin{figure}[tb]\n\\centerline{\\includegraphics[angle=0,width=4in]{H024_4.smaller.ps}}\n\\vspace*{3mm}\n\\caption{Snapshot from PIMC simulations of pure hydrogen in the molecular regime at \n $T=5\\,000\\,\\rm K$ and $r_s=4.0$. The pink spheres denote the\n protons. The bonds (white lines) were added as a guide to the\n eye. The electron paths are shown in red and blue [light and\n dark gray] depending on their spin state. }\n\\label{pict_dilute_H2}\n\\end{figure}\n\nAt low density and temperature, one finds a fluid of interacting\nhydrogen molecules. A PIMC snapshot for $T=5\\,000\\,\\rm K$ and\n$r_s=4.0$ is shown in Fig.~\\ref{pict_dilute_H2}. The proton paths are\nvery localized due to their high mass. Their spread can be estimated\nfrom the de Broglie thermal wavelength, $\\lambda_{d}^{2}=\\hbar^{2}\n\\beta \/2\\pi m$. The electron paths are more spread out but they \nare localized to some extent since two electrons of opposite spin\nestablish the chemical bond in the hydrogen molecule.\n\nIf the temperature is raised from 5000 to $10^6$ K, hydrogen\nundergoes a smooth transition from a molecular fluid though an atomic\nregime and finally to a two-component plasma of interacting electrons\nand protons. Many-body simulations at even higher temperatures are not\nneeded since analytical methods like the Debye-H{\\\"u}ckel theory work\nwell. PIMC with explicit treatment of the electrons can also be\napplied to temperatures below 5$\\,$000$\\,$K but groundstate methods\nare then more practical since excitation become less relevant and the\ncomputational cost scale with the length of paths like $T^{-1}$.\n\nA detailed analysis of the chemical species present in the low density\nregime ($r_s \\ge 2.6$) is given in Ref.~\\onlinecite{MC00}. With\ndecreasing density, one finds that the degree of molecular\ndissociation increases since the atomic state has higher entropy. For\nthe same reason, one observes that the degree of atomic ionization\nincreases with decreasing density. All these low-density effects can\nbe well characterized by analytical models based on approximate free\nenergy expressions for atoms, molecules and ionized\nparticles.\\cite{SC92,Ro90,rotesbuch,Ju01} In this regime, one also\nfinds reasonably good agreement between for the EOS derived from PIMC\nand chemical models.\\cite{MC00}\n\n\\begin{figure}[tb]\n\\centerline{\\includegraphics[angle=0,width=4in]{D023_1.smaller.ps}}\n\\caption{Deuterium in the dense molecular regime ($T=5\\,000\\,\\rm K$ and $r_S=1.86$) \n Due to the density increase compared to\n Fig.~\\ref{pict_dilute_H2} (see details there), the electron\n paths permute with a rising probability (shown as yellow [light\n gray] lines) but are still localized enough to form a bond\n between the two protons in the molecule. The electron density\n average over many electron configurations is indicated in gray\n color on the blue rectangles.}\n\\label{pict_dense_H2}\n\\end{figure}\n\nIf the density is increased at $T=5\\,000\\,$K, one finds an\nintermediate regime of strongly interacting molecules\n(Fig.~\\ref{pict_dense_H2}). Some electron paths exchange with\nneighboring molecules indicating the importance of fermionic\neffects. However, the electrons are still localized enough to provide\na sufficient binding force for the protons.\n\n\\begin{figure}[tb]\n\\centerline{\\includegraphics[angle=0,width=4in]{D159_1.smaller.ps}}\n\\caption{Deuterium in the metallic regime characterized by unpaired \n\tprotons and a gas of degenerate electrons. The snapshot was\n\ttaken from a PIMC simulation at $T=6\\,250\\,\\rm K$ and\n\t$r_s=1.60$. The electron paths are delocalized and exchange\n\tfrequently (see description of Fig.~\\ref{pict_dense_H2}).}\n\\label{pict_metallic_H}\n\\end{figure}\n\nIf the density is increased further from $r_s=1.86$ to 1.60, this\nbinding force is lost due to further delocalization of the electrons\n(Fig.~\\ref{pict_metallic_H}). Almost all electrons are now involved in\nlong exchange cycles indicating a highly conducting, metallic\nstate. No binding forces of the protons can be observed. One can\ntherefore conclude that hydrogen metallizes in dissociated, or atomic\nform, a conclusion that is consistent with DFT\nsimulations\\cite{Banier2000,Scandolo2003,Desjarlais2003,Bonev2004}. If\nthe density is increased further, the electrons form a rigid\nneutralizing background and one recovers the limit of a one-component\nplasma of protons. If the temperature is increased, the electron paths\nget shorter and shorter, the degree of electron degeneracy decreases\ngradually and one recovers the limit of a two-component plasma at high\ntemperature.\n\nThe nature of the molecular-metallic transition has not yet been\ndetermined with certainty. A large number of analytical\nmodels\\cite{Be99} predict a first order transition, others do\nnot.\\cite{Ro90} PIMC simulation with free particle nodes by\nW. Magro~{\\em et al.}\\cite{Ma96} showed an abrupt transition characterized by\na region $\\left.\\frac{dP}{dT}\\right|_V<0$ for $r_s$=1.86. Later work\nby B. Militzer\\cite{BM00} using more accurate variational nodes did\nnot show such a region for $r_s$=1.86. Whether PIMC simulations with\nvariational nodes predict a gradual molecular-metallic transition for\nall temperatures, or if the region of an abrupt transition has shifted\nto lower, not yet accessible, temperatures has not been\ndetermined. However, recent density functional molecular dynamics\nsimulations by various authors predict a first order transition below\n$5\\,000\\,$K.\\cite{Banier2000,Scandolo2003,Desjarlais2003,Bonev2004}\n\nThere was a lot of recent interest focused on the hydrogen EOS because\nof the {\\em unexpectedly high compressibility} inferred from the\nlaser-driven shock wave experiment by Da Silva~{\\em et al.}\\cite{Si97} and\nCollins~{\\em et al.}\\cite{Co98} using the Nova laser at Lawrence Livermore\nNational Laboratory. The measurements indicated that hydrogen could be\ncompressed to about 6-fold the initial density rather than 4-fold as\nindicated by the Sesame model.\\cite{Ke83,Kerley2003} Even though the\ntemperatures and pressures reached in shock experiments are higher\nthan those in the giant planets these measurements are crucial here\nsince they represent the only way to distinguish between different EOS\nmodels at high temperature.\n\nThe Nova results challenged the existing understanding of high P-T\nhydrogen and triggered many new experimental and theoretical\nefforts. Different chemical models gave rise to very different\npredictions\\cite{Be99,Ro90,Ju00} ranging from 4-fold compression as\nsuggested in\\cite{Ke83,Kerley2003} to 6-fold as predicted by the Ross\nmodel.\\cite{Ro98} While the accuracy of chemical models did not allow\nany conclusive predictions, first principle simulations from\nPIMC\\cite{MC00,Mi01} as well as from density functional molecular\ndynamics\\cite{Le00,Desjarlais2003,Bonev2004} consistently predict a\nlower compressibility of about 4.3.\n\nSince then there have been many attempts to resolve this discrepancy\nbut the most important contributions came from new experiments by\nKnudson~{\\em et al.}\\cite{Kn01,Kn03,Kn04} at Sandia National\nLaboratory. Instead of a laser drive, they used magnetically driven\nshock waves in combination with bigger samples. They found a\nsignificantly lower compressibility quite similar to predictions from\nfirst principles methods. The new results are also supported by a\nthird set of experiments by Russian investigators using spherically\nconverging shock waves.\\cite{Belov2002,Boriskov2003}\n\n\\section{HYDROGEN-HELIUM MIXTURES}\n\nDue to importance for astrophysical applications, EOS models for\nhy\\-dro\\-gen-he\\-lium mixtures have been studied for quite some time. The\nmost widely used EOS was derived by Saumon, Chabrier and van\nHorn (SCH).\\cite{SC95} Like previous chemical models, it is based on a\nhydrogen and a helium EOS that combined using an ideal mixing\nrule. \n\nFollowing the discussion of the hydrogen EOS above, we are now\nanalyzing the second ingredient: the EOS of helium. Given its low\nabundance in giant planets, helium is not expected to be present in\nvery high concentration. However,\nStevenson\\cite{Stevenson77b} has proposed\nthat the hydrogen-helium mixture could phase separate under certain\nhigh pressure conditions. As a result, helium droplets would form and\nfall as rain towards planet core. The associated release of latent\nheat would delay the cooling process of the planet. This was suggested\nas one possible explanation of why standard models predict Saturn to cool\nat a much faster rate than is observed.\\cite{FH04} In order to\ndetect hydrogen-helium phase separation one needs an EOS of pure\nhelium that we will now discuss.\n\n\\begin{figure}[htb]\n\\centerline{\\includegraphics[width=2.48in]{He_P_vs_T06.eps}~~~~\\includegraphics[width=2.4in]{He_E_vs_T08.eps}}\n\\caption{Equation of state comparison for pure helium showing the pressure (left) \n and the internal energy (right) for different density as\n derived from path integral Monte Carlo simulations and the\n chemical model by Saumon, Chabrier, and van Horn.\\cite{SC95}}\n\\label{HeEOS}\n\\end{figure}\n\nFigure~\\ref{HeEOS} shows an EOS comparison of our PIMC calculations\nwith the SCH model. One finds that the pressure is underestimated by\nSCH and that the deviations increase with density. The internal energy\nat high temperature is also underestimated. The analysis suggests a\nmore careful treatment of the helium ionization states is needed to\nimprove the chemical model.\n\nTo conclude the comparison with chemical models, we test one more\nassumption that is generally made: the ideal mixing hypothesis. We\nperformed a large number of PIMC calculations of fully interacting\nhydrogen-helium mixtures at a fixed density of $r_s=1.86$ for various\ntemperatures and mixing ratios. The resulting correction to ideal\nmixing is given by,\n\\begin{equation}\n\\Delta f_{\\rm mix} = f(x) - (1-x)\\,f(x=0) - x\\,f(x=1)\\,.\n\\end{equation} \n\n\\begin{figure}[htb]\n\\centerline{\\includegraphics[width=2.48in]{mixing_pressure02.eps}~~~~\\includegraphics[width=2.4in]{mixing_internal_energy02.eps}}\n\\caption{Excess mixing pressure $\\Delta P_{\\rm mix}$ and internal energy per electron $\\Delta E_{\\rm mix}$\n are shown for a hydrogen-helium mixture at $r_s=1.86$. Path\n integral Monte Carlo results are shown for several\n temperatures. The mixing is performed at constant\n density. $\\Delta E_{\\rm mix}$ is largest for T=31$\\,$250$\\,$K\n because at this temperature the two end members are\n characterized by very different degrees of ionization.}\n\\label{mix}\n\\end{figure}\n\nFigure~\\ref{mix} shows that the corrections to pressure, $\\Delta\nP_{\\rm mix}$, increase with decreasing temperature reaching 10\\% for\n$15\\,625\\,$K. In the considered temperature interval, the corrections\nto the internal energy are largest for $31\\,250\\,$K which can be\nexplained by the different ionization states of the two fluids that\nlead to larger error if ideal mixing is assumed.\n\n\\begin{figure}[htb]\n\\centerline{\\includegraphics[width=4in]{gr05.eps}}\n\\caption{Comparison of different pair correlation functions, $g(r)$, from three PIMC \n simulations at $r_s=1.86$ and $15\\,625\\,$K: pure hydrogen (red\n dotted lines), a $x=50\\%$ mixture (black dashed lines) and pure\n helium (blue solid lines). The six graphs show $g(r)$ for\n different pairs of protons (p), electrons (e) and helium nuclei\n ($\\alpha$). For electrons, pairs with antiparallel (upper\n right) and with parallel (lower right) spins are\n distinguished.}\n\\label{fig:gr}\n\\end{figure}\n\nFinally, we discuss pair correlation functions, $g(r)$, for hydrogen-helium mixtures, \n\\begin{equation}\n g(r)\\equiv \\frac{\\Omega}{N^2}\n \\left<\\sum_{i\\neq j}\\delta ({\\bf r}-{\\bf r}_{ij})\n \\right>\\;.\n \\label{eq3.1}\n\\end{equation}\n$g(r)$ functions are a standard tool to characterize the short-range\ncorrelation of particles in liquids. In Fig.~\\ref{fig:gr}, we compare\nthe correlation functions for pure hydrogen, a $x$=50\\% mixture, and\npure helium at fixed temperature (15$\\,$625$\\,$K) and density\n($r_s=1.86$). Changes in the pair correlations are now discussed as a\nfunction of helium concentration. The electron-proton $g(r)$ shows a\nstrong peak at the origin due to the Coulomb attraction. The peak is\nnot affected when 50\\% helium is added to a pure hydrogen\nsample. Similarly, the correlation between electrons and helium nuclei\nis not altered by the presence of protons.\n\nThe electron-electron correlation depends strongly on spin. For pairs\nwith parallel spin, exchange effects lead to a strong repulsion. The\nresulting $g(r)$ function does not depend much on whether helium is\npresent or not. The correlation function between pairs of antiparallel\nspins, on the other hand, is strongly affected by the presence of\nhelium nuclei. In the helium atom, two electrons with opposite spin\nare attracted to the core, which indirectly leads to the observed\nincrease in the electron-electron $g(r)$.\n\nIt is interesting to note that proton-proton $g(r)$ changes with the\nhelium concentration but the correlation of helium nuclei does\nnot. With increasing presence of helium, a peak in the proton-proton\n$g(r)$ appears at $r=1.4\\,a_0$ which indicates the formation of H$_2$\nmolecules. Adding helium nuclei leads to the localization of a\nfraction of the electrons. The available space in combination with the\nreduced electronic exchange effects then leads to the formation of hydrogen\nmolecules which are not present in pure hydrogen at the same\ntemperature and density.\n\n\\section{CONCLUSIONS}\n\nThe discussion in this article centered on the phase diagram of hot\ndense hydrogen. The first PIMC results for hydrogen-helium mixtures\nwere presented here, and the accuracy of existing helium EOS models as\nwell as the commonly used ideal mixing rule were analyzed. This\nanalysis revealed substantial inaccuracies in existing EOS\nmodels. These uncertainties prevent us from making reliable models for\nthe interior of giant planets and from drawing conclusions about their\nformation mechanism, either nucleated capture of nebular\nhydrogen\\cite{Li93} or gravitational instability driven formation in\nprotostellar disks.\\cite{Boss97,Boss04}\n\nAn improved EOS will also yield a better characterization of the 120\nextrasolar giant planets that have been discovered so far using radio\nvelocity measurements.\\cite{Ma02} In particular the mass-radius\nrelationship obtained from transit extrasolar giant\nplanets\\cite{brown01} will allow one to draw conclusions about the\nheavy element abundance in those planets. The metallicity of a planet\nrelative to its central star constrains the role of accretion of\nplanetesimals in the planet's formation.\n\n\\section*{ACKNOWLEDGMENTS}\n\nThe author thanks W. Hubbard and D. Stevenson for useful discussions\nand D. Saumon for providing us with the SCH EOS.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nThe mass distribution of neutron stars contains information about the\nsupernova explosion mechanisms, the equation of state of neutron star\nmatter, and the accretion history of each neutron star since its\nformation. Certain populations of neutron stars such as those in\ndouble neutron stars and in binaries with high mass companions are\nthought to have experienced little-to-no accretion over their\nlifetimes. In contrast, neutron stars in low-mass X-ray binaries and\nfast pulsars, which are typically in close orbits around white dwarfs,\nundergo extended accretion periods that are likely to move the neutron\nstar mass away from its birth value.\n\nThe neutron star mass measurements that were available a decade ago\nallowed a statistical inference of the mass distribution of double\nneutron stars (Finn 1994) or of pulsars in binaries, without\ndistinguishing between subgroups (Thorsett \\& Chakrabarty 1999). Finn\n(1994) found that neutron star masses fall predominantly in the\n$1.3-1.6~M_\\odot$ range. Thorsett \\& Chakrabarty (1999) found that the\nmass distribution for the combined population is consistent with a\nnarrow Gaussian at $1.35 \\pm 0.04~M_\\odot$. More recently, Schwab,\nPodsiadlowski, \\& Rappaport (2010) argued that the distribution of\nneutron star masses in double neutron stars is actually bimodal, with\none peak centered at $\\sim 1.25~M_\\odot$ and the other at $\\sim\n1.35~M_\\odot$, which they attributed to different supernova explosion\nmechanisms. Kiziltan, Kottas, \\& Thorsett (2010), Valentim, Rangel, \\&\nHorvath (2011), and Zhang et al.\\ (2011), on the other hand, inferred\nthe mass distribution of different neutron star subgroups based either\non the pulsar spin period or the binary companion, both of which were\ntaken to be indicative of the accretion history of the system. All\ngroups found that the neutron stars that are thought to have undergone\nsignificant accretion are on average $0.2-0.3~M_\\odot$ heavier than\nthose that have not.\n\nOne result that is common to all of these studies is the narrowness of\nthe mass distribution of double neutron stars, $\\sigma \\simeq\n0.05~M_\\odot$, which has been taken as indicative of the birth mass\ndistribution of all neutron stars. The mean of the distribution is at\n$1.35~M_\\odot$, which is significantly larger than the mass of the\npre-supernova iron core for neutron stars that form through the\ncore-collapse mechanism. The Chandrasekhar mass for cores with\nelectron fractions in the range $Y_e = 0.42-0.48$ is\n$1.15-1.34~M_\\odot$. Electrostatic interactions and entropy of the\ncore introduce additional corrections to the pre-collapse mass (see\nTimmes et al.\\ 1996 for a discussion). Taking into account the\nbinding energy of the neutron star results in gravitational masses for\nthe collapsed cores in the range $1.06-1.22~M_\\odot$. Even the largest\nof these masses is well below the mean of the observed distribution of\ndouble neutron stars. Fallback of stellar matter onto the collapsing\ncore during the supernova explosion allows for the remnant to\nincrease. However, this is also expected to increase the dispersion of\nmasses by a comparable amount (see Zhang et al. 2008), which is\ninconsistent with the narrowness of the inferred mass distribution of\ndouble neutron star masses.\n\nConsidering a bimodal underlying distribution in the population of\ndouble neutron stars, as in Schwab et al.\\ (2010), makes the width of\neach distribution even narrower: $0.008~M_\\odot$ and $0.025~M_\\odot$\nfor the two components. For the lower mass component centered around\n$\\sim 1.25~M_\\odot$, such a narrow distribution may be reasonably\nobtained through an electron capture supernova, the onset of which\noccurs at a particular mass threshold of an ONeMg white dwarf\n(Podsiadlowski et al.\\ 2005). However, the second component, which is\ncentered at $1.35~M_\\odot$ cannot be explained as a result of the\nelectron capture supernovae and poses the same challenge in its\nnarrowness when explained via the core-collapse mechanism.\n\nIn order to model the distribution of neutron star masses both at\ntheir births and throughout their lives, one important question to\naddress is whether double neutron stars are a representative sample\nfor neutron stars at their birth masses. In this paper, we address\nthis question by identifying a different population of neutron stars\nat or near their birth masses and compare the inferred mass\ndistribution with that of double neutron stars. Furthermore, to\npinpoint the effects of subsequent accretion, we compare the inferred\nmass distribution of these neutron stars to that of neutron stars\nwhich have undergone or are currently undergoing extended mass\naccretion.\n\nMaking use of all of the currently available neutron star mass\nmeasurements, we divide the sample into various subgroups based on the\nnature of the companion as well as the neutron star spin. We employ a\nuniform Bayesian statistical approach that utilizes the entire\nposterior likelihood of each mass measurement to infer the parameters\nof the underlying mass distribution without assuming Gaussian errors.\n\nIn Section 2, we present the neutron star mass measurements, grouping\nthe data according to the measurement technique and the type of the\nbinary. In Section 3, we estimate the parameters of the underlying\nmass distribution for each subgroup and assess the sensitivity of our\nresults on the particular measurements and priors. In section 4, we\nsummarize our findings and discuss their implications for the\ndifferent physical mechanisms that determine the neutron star mass\ndistribution. \n\n\\section{Mass Measurements}\n\nMass measurements of neutron stars are carried out in several\ndifferent ways. When neutron stars appear as pulsars, high precision\npulse timing observations lead to a measurement of several orbital\nparameters. The pulsar's orbit can be described in classical gravity\nby the five Keplerian parameters: the binary period $P_b$, the\neccentricity $e$, the projection of the pulsar's semi-major axis on\nthe observer's line of sight $a \\sin i$, the time of periastron $T_0$,\nand the longitude of periastron $\\omega_0$, where $i$ is the angle\nbetween the orbital angular momentum vector and the line of sight. The\nmass function, which is related to the mass of the pulsar $M_{\\rm psr}$,\nits companion $M_{\\rm c}$, and the inclination angle $i$, \n\\begin{equation}\nf = \\frac{(M_{\\rm c} \\sin i)^3}{M_{\\rm T}^2} = \\left(\\frac{2 \\pi}\n{P_b}\\right)^2 \\frac{(a \\sin i)^3}{G}, \n\\end{equation}\nis, therefore, directly obtained from these orbital parameters, where\n$M_{\\rm T} = M_{\\rm psr} + M_{\\rm c}$ is the total mass of the system.\n\nProceeding from a mass function to a measurement of the mass of the\npulsar and of its companion requires additional information. This\ninformation can come from measurement of relativistic effects in the\nbinary orbits or from independent observations of the companion\nstars. We discuss below the various techniques, the measurements they\nresulted in to date, and the associated uncertainties.\n\nThe majority of the precise neutron star mass measurements come from\nradio pulsar timing techniques and rely on the measurement of\nrelativistic effects in the binary orbits.\n\nThe general relativistic effects can be described by five additional\n``post-Keplerian'' (or PK) parameters. These are: the advance of\nperiastron $\\dot{\\omega}$, the orbital period decay $\\dot{P}_{b}$, the\ntime dilation-gravitational redshift factor $\\gamma$, as well as the\nrange $r$ and the shape $s$ of Shapiro delay, which are related to the\ncomponent masses, the orbital period, and eccentricity by\n\\begin{equation}\n\\dot{\\omega} = 3\\left(\\frac{P_{b}}{2\\pi}\\right)^{-5\/3}\n\\left( \\frac{G M_{\\rm T}}{c^3} \\right)^{2\/3}\\left(1-e^2\\right)^{-1},\n\\end{equation}\n\\begin{equation}\n\\dot{P}_b = -\\frac{192\\pi}{5}\\left(\\frac{P_{b}}{2\\pi}\\right)^{-5\/3} \n\\left(\\frac{G}{c^3}\\right)^{5\/3} \n\\left(1+\\frac{73}{24}e^2+\\frac{37}{96}e^4\\right)\\times \\nonumber\\\\\n(1-e^2)^{-7\/2}\\, M_{\\rm psr} M_{\\rm c}\\,M_{\\rm T}^{-1\/3} ,\n\\end{equation}\n\\begin{equation}\n\\gamma = e\\left(\\frac{P_{b}}{2\\pi}\\right)^{1\/3}\n\\left(\\frac{G}{c^3}\\right)^{2\/3} M_{\\rm T}^{-4\/3}\nM_{\\rm c}\\left(M_{\\rm psr} + 2 M_{\\rm c} \\right) ,\n\\end{equation}\n\\begin{equation}\nr = \\frac{G}{c^3} M_{\\rm c},\n\\end{equation}\n\\begin{equation}\ns = G^{-1\/3} a\\sin i \\left(\\frac{P_{b}}{2\\pi}\\right)^{-2\/3}\nM_{\\rm T}^{2\/3}\\, M_{\\rm c}^{-1}.\n\\label{Eq:PK}\n\\end{equation}\n\nIn highly eccentric systems that have been observed repeatedly over a\nlong period of time, the measurement of $\\dot{\\omega}$ is usually\npossible and leads to a strong constraint on the total mass of the\nbinary $M_{\\rm T}$. The measurement of the parameter $\\gamma$ requires\nsimilarly eccentric systems and long-term monitoring. In high\ninclination systems, on the other hand, it is sometimes possible to\ndetect the Shapiro time delay and obtain the parameters $r$ and\n$s$. The measurement of the rate of orbital period decay typically\nrequires the longest monitoring and timing of the pulsar, sometimes\nover decades. This has been achieved for a handful of pulsars in\nbinaries.\n\nThe precision with which the pulsar mass can be determined ultimately\ndepends on the number of PK parameters that are measured for that\nbinary. In systems where two or more PK parameters are known, the\npulsar mass is precisely determined. In the category where only one PK\nparameter is known in addition to the mass function, the mass of each\nneutron star is not as well constrained.\n\nIn the following, we divide neutron star mass measurements into\ncategories based on the information available for each binary system\nsuch as the number of PK parameters. For each category, we derive the\nlikelihood $P_i({\\rm data} \\vert M_{\\rm NS})$, which measures the\nchance of obtaining the particular set of data for the $i$-th source\nif that source had mass $M_{\\rm NS}$.\n\nWe are ultimately interested in delineating the effects of mass\naccretion from the neutron star birth masses. We, therefore, further\ndivide each category into groups based on the nature of the companion\nstar or the spin of the pulsar. In particular, neutron star-neutron\nstar (NS-NS) binaries, as well as eclipsing X-ray pulsars in high-mass\nX-ray binaries and slow radio pulsars are expected to have experienced\nlittle-to-no accretion. On the other hand, neutron star-white dwarf\n(NS-WD) binaries are the remnants of a long-lasting low-mass X-ray\nbinary phase, where significant mass accretion may have occurred. We\nalso group millisecond pulsars with main sequence companions (NS-MS)\nalong with the latter group, because of the probable recycling these\nneutron stars underwent to reach millisecond periods. We will refer to\nthe latter group as ``fast pulsars''. Finally, we will also consider\naccreting X-ray bursters, for which the masses have been measured\nprimarily through X-ray spectroscopy.\n\nIn detail, the various categories are:\n\n{\\it (Ia) Double Neutron Stars with at least 2 PK Parameters.---\\\/}\nSix neutron star systems shown in Table~1 have at least 2 measured PK\nparameters leading to well determined masses. In this case, the\nlikelihood of neutron star masses is highly symmetric and narrowly\npeaked, and thus can be described as a Gaussian\n\\begin{equation}\nP_i({\\rm data} \\vert M_{\\rm NS}) = C_i \\exp\\left[-\\frac{(M_{\\rm NS}-M_{0,i})^2}\n{2\\sigma_{M,i}^2} \\right]\n\\end{equation}\nwith a mean $M_{0,i}$ and a standard deviation $\\sigma_{M,i}$. In this\nand the following expressions, $C_i$ is a proper normalization\nconstant such that\n\\begin{equation}\n\\int_0^\\infty P_i({\\rm data} \\vert M_{\\rm NS}) dM_{\\rm NS} = 1.\n\\end{equation} \n\n\\input{table1.tex}\n\nWe plot in Figure~\\ref{fig:psr_gauss} the likelihood of the masses\nof NS-NS binaries that belong to this category.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f1.ps}\n\\caption{The likelihoods $P_i({\\rm data}\\vert M_{\\rm NS})$ for the 12\ndouble neutron stars with precisely determined masses arising from the\nmeasurement of $\\ge 2$ PK parameters. These systems belong to category\nIa discussed in the text.}\n\\mbox{}\n\\label{fig:psr_gauss} \n\\end{figure}\n\n\n{\\it (Ib) Neutron Stars in Binaries with at least 2 PK\nParameters.---\\\/} To date, observations of nine neutron stars that are\nin close orbits around predominantly white dwarf companions have\nyielded a measurement of at least 2 PK parameters. The known exception\nis PSR~J1903+0327, with a probable main sequence companion, but is\nnevertheless thought to be recycled owing to its millisecond\nperiod. As in the previous category, we assign a Gaussian likelihood\nto each measurement with a mean $M_{0,i}$ and a standard deviation\n$\\sigma_{M,i}$. We present in Table~2 and\nFigure~\\ref{fig:psr_recycl_gauss} the mass measurements and their\nuncertainties for these NS-WD binaries. Even though it has a white\ndwarf companion, PSR J1141$-$6545 is different from the rest of the\nsources in this category in that it is a slowly spinning neutron\nstar. For this reason, we group it with the accreting and slow pulsars\ndiscussed below and show its likelihood in Figure~\\ref{fig:psr_slow}.\n\n\\input{table2.tex}\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f2.ps}\n\\caption{The likelihoods $P_i({\\rm data}\\vert M_{\\rm psr})$ for nine\nrecycled pulsars with white dwarf companions. The measurement of $\\ge\n2$ PK parameters in these binaries lead to precisely determined\nneutron star masses. These systems belong to category Ib discussed in\nthe text.}\n\\mbox{}\n\\label{fig:psr_recycl_gauss} \n\\end{figure}\n\n{\\it (IIa) Double Neutron Stars with 1 PK Parameter.---\\\/} \nIn this category, there are three double neutron stars, for which \nthe measurement of the advance of periastron $\\dot{\\omega}$\nallows for a precise determination of the total mass of the binary.\nHowever, in the absence of a second PK parameter or a knowledge of the\nsystem inclination $i$, the mass of each neutron star is not well\nconstrained. In these cases, using the total mass of the system as \na constraint allows us to write the likelihood of the mass of the \npulsar as\n\\begin{equation} \nP_i({\\rm data} \\vert M_{\\rm psr}) = C_i \\int dM_{\\rm tot} \n\\exp\\left[-\\frac{(M_{\\rm tot}-M_{{\\rm tot},0})^2}\n{2\\sigma_{M_{\\rm tot},i}^2} \\right] \\int d(\\cos i)\n\\exp\\left[-\\frac{\\left(f_0-\\frac{M^3_{\\rm psr}\\sin^3i}{M^2_{\\rm tot}}\\right)^2}\n{2\\sigma_f^2}\\right]\n\\label{eq:pm_omegadot}\n\\end{equation}\nand of the companion as \n\\begin{equation} \nP_i({\\rm data} \\vert M_{\\rm c}) = C_i \\int dM_{\\rm tot} \n\\exp\\left[-\\frac{(M_{\\rm tot}-M_{{\\rm tot},0})^2}\n{2\\sigma_{M_{\\rm tot},i}^2} \\right] \\int d(\\cos i)\n\\exp\\left[-\\frac{\\left(f_0-\\frac{(M_{\\rm tot}-M_{\\rm c})^3\n\\sin^3i}{M^2_{\\rm tot}}\\right)^2}\n{2\\sigma_f^2}\\right]. \n\\label{eq:pm__c_omegadot}\n\\end{equation}\n\nWe present in Table~3 the relevant pulsar data for the three double\nneutron star binaries that are in this category.\nFigure~\\ref{fig:psr_omega} shows the likelihood of each neutron star\nmass $P_i({\\rm data}\\vert M_{\\rm NS})$ for these sources, with the top\npanel including the pulsars and the bottom panel the companions. Note,\nhowever, that even though they are shown independently for the\npurposes of this figure, the likelihoods of the masses of the pulsar\nand its companion are not independent probabilities. Therefore, when\ninferring the mass distribution of double neutron stars, the\nconstraint over the total mass is incorporated as we will discuss in\n\\S3.\n\n\\input{table3.tex}\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f3a.ps}\n \\includegraphics[scale=0.75]{f3b.ps}\n\\caption{The likelihoods $P_i({\\rm data}\\vert M_{\\rm NS})$ for the\ndouble neutron stars with 1 PK parameter for (top) the pulsars and\n(bottom) companion neutron stars. These belong to category IIa\ndiscussed in the text.}\n\\mbox{}\n\\label{fig:psr_omega} \n\\end{figure}\n\n{\\it (IIb) Neutron Stars in Binaries with 1 PK Parameter.---\\\/} This\ncategory comprises of eleven neutron star binaries with mostly white\ndwarf companions. In several systems, the companions have not been\nidentified. A constraint on the total mass comes from the measurement\nof the rate of advance of the periastron. We assign to each neutron\nstar mass a likelihood according to equation~(\\ref{eq:pm_omegadot}).\nTable~4 and Figure~\\ref{fig:psr_recycl_omega} shows the relevant\nparameters for these systems. Note that, of this category, PSR B2303+46\nis a slowly spinning neutron star. For this reason, we group it with\nthe accreting and slow pulsars discussed below and show its likelihood\nin Figure~\\ref{fig:psr_slow}.\n\n\\input{table4.tex}\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f4.ps}\n\\caption{The likelihoods $P_i({\\rm data}\\vert M)$ for \nthe recycled neutron stars with 1 PK parameter (blue curves; category\nIIb) or with optical observations of the white dwarf companions (red\ncurves; category III). }\n\\mbox{}\n\\label{fig:psr_recycl_omega} \n\\end{figure}\n\n{\\it (III) Neutron Stars in Binaries with Optical Observations of\nWhite Dwarf Companions.---\\\/} For two neutron stars in orbit around\nwhite dwarfs, optical observations of the companions have resulted in\nthe measurements of the mass $M_{\\rm WD}$ as well as of the radial\nvelocity amplitude $K_{\\rm opt}$ of the white dwarf. The latter, in\ncombination with the orbital parameters obtained from the radio timing\nsolution gives the mass ratio of the binary according to\n\\begin{equation}\nq = K_{\\rm opt} P_{\\rm b} \\frac{(1-e^2)^{1\/2}}{2 \\pi a_{\\rm psr} \\sin i}.\n\\end{equation}\nUsing this information, we can then calculate the likelihood of the \nneutron star mass as \n\\begin{equation}\nP_i({\\rm data} \\vert M_{\\rm psr}) = C_i \\int dM_{\\rm WD} \\int di \\sin i\n\\exp\\left[-\\frac{(M_{\\rm WD}-M_{{\\rm WD},0})^2}\n{2\\sigma_{M_{\\rm WD},i}^2} \\right] \n\\exp\\left[-\\frac{(\\frac{M_{\\rm psr}}{M_{\\rm WD}}-q_{i,0})^2}\n{2\\sigma_{q,i}^2} \\right] \\delta[f(i)-f_{0,i}],\n\\end{equation}\nwhere we assume that the error in the measurement of each mass function \nis negligible. We perform the integration over inclination making use of \nthe identity\n\\begin{equation}\n\\delta[f(i)-f_{0,i}] = \\frac{\\delta(i-i_0)}{\\vert df\/di\\vert_{i_0}}\n\\end{equation}\nwhere $i_0$ is the solution of the equation $f(i_0) - f_{0,i}=0$, i.e.,\n\\begin{equation}\n\\sin i_0 = \\left[f_{0,i} \\left(1+\\frac{M_{\\rm psr}}{M_{\\rm WD}}\\right)^2 \nM_{\\rm WD}^{-1}\\right]^{1\/3}.\n\\end{equation}\nGiven that \n\\begin{equation}\n\\left\\vert \\frac{df}{di}\\right\\vert_{i_0} = 3 f \\left\\vert \n\\frac{\\cos i_0}{\\sin i_0}\\right\\vert, \n\\end{equation}\nthe likelihood becomes\n\\begin{equation}\nP_i({\\rm data} \\vert M_{\\rm psr}) = \\frac{C_i}{3 f_{0,i}} \n\\int dM_{\\rm WD} \\frac{\\sin^2 i_0}{\\cos i_0}\n\\exp\\left[-\\frac{(M_{\\rm WD}-M_{{\\rm WD},0})^2}\n{2\\sigma_{M_{\\rm WD},i}^2} \\right]\n\\exp\\left[-\\frac{(\\frac{M_{\\rm psr}}{M_{\\rm WD}}-q_{i,0})^2}\n{2\\sigma_{q,i}^2} \\right]\\;.\n\\end{equation}\n\nTable~5 summarizes the relevant data for these two binaries. The\nlikelihoods of the neutron star masses are shown as red curves in\nFigure~\\ref{fig:psr_recycl_omega}.\n\n\\input{table5.tex}\n\n{\\it (IV) Eclipsing X-ray Pulsars.---\\\/} Eclipsing X-ray pulsars\nprovide a wealth of observational information using which we can\nestimate the masses of these neutron stars. X-ray observations of\neach pulsar give the orbital period of the binary $P_b$, the\neccentricity of the orbit $e$, longitude of periastron $\\omega_0$, the\nsemi-major axis of the neutron star's orbit $a_X\\sin i$, and the\nsemi-duration of the eclipse $\\theta_e$. In addition, optical\nobservations of the companion star give its velocity amplitude $K_{\\rm\nopt}$, its projected rotational velocity $v_{\\rm rot} \\sin i$, and the\namplitude of ellipsoidal variations $A$. From these observables, it is\npossible to solve for the fundamental parameters of the binary, viz.,\nthe mass of the neutron star $M_{\\rm NS}$, the mass, radius and\nrotational angular velocity of the companion, $M_{\\rm opt}$, $R_{\\rm\nopt}$, $\\Omega_{\\rm opt}$, and the inclination angle of the binary\n$i$.\n\nThe necessary data are available for six eclipsing pulsars: Vela X--1,\n4U1538--52, SMC X--1, LMC X--4, Cen X--3, and Her X--1 (see, e.g., van\nKerkwijk, van Paradijs, \\& Zuiderwijk 1995; van der Meer et\nal. 2007). Recently, Rawls et al. (2011) collected all the available\ndata and presented a detailed analysis of the likelihood of mass for\nthe individual neutron stars; see Table~6 for a compilation of the\nrelevant results. \n\nThere are several sources of potential systematic uncertainties in the\nmasses inferred in these eclipsing binaries. For example, there are\nsignificant residuals in the radial velocity curves of Vela X-1 after\nthe best-fit orbital solution is subtracted (e.g., Barziv et\nal. 2001), in which Quaintrell et al. (2003) noted the presence of a\nperiodicity and suggested modes on the star as a possible origin. More\nrecently, Koenigsberger et al.\\ (2012) developed a model for these\nresiduals based on the interaction between the neutron star and its\ncompanion but did not directly fit the model to the data. Even though\nthey concluded that a 1.55~$M_\\odot$ neutron star is marginally\nconsistent with observations, a higher mass for the neutron star,\ncentered around 1.7~$M_\\odot$, appears to be favored.\n\nA second important source of systematic uncertainty arises from\nmodeling of the ellipsoidal modulations, which includes a contribution\nfrom the accretion disk and may significantly influence the inferred\nneutron star masses. This disk contribution is often non-negligible\n(see Figs.\\ 8, 9 of Rawls et al.\\ 2011) and involves multiple\nparameters. In order to assess the influence of this additional\ninformation on our results, we reanalyze here the data of eclipsing\npulsars ignoring the ellipsoidal modulations.\n\nFor each of the six systems, we write the likelihood of the data as a\nfunction of the neutron star mass $M_{\\rm NS}$ as follows,\n\\begin{eqnarray}\nP(\\rm{data}|M_{\\rm NS}) &=& C \\int_{M_{\\rm opt,min}}^{M_{\\rm opt,max}} dM_{\\rm opt}\n\\int_{\\beta_{\\rm min}}^{\\beta_{\\rm max}} d\\beta\n\\int_{\\Omega_{\\rm opt,min}}^{\\Omega_{\\rm opt,max}} d\\Omega_{\\rm opt}\n\\int_{0}^{1} d(\\cos i) \\nonumber \\\\\n&~& ~~~\\exp\\left[-\\frac{(f_M-f_0)^2}{2\\sigma_f^2}\n-\\frac{(K_{\\rm opt}-K_0)^2}{2\\sigma_K^2}\n-\\frac{(v_{\\rm rot}\\sin i-v_0)^2}{2\\sigma_v^2}\n-\\frac{(\\theta_e-\\theta_0)^2}{2\\sigma_\\theta^2}\n\\right].\n\\label{eclipse}\n\\end{eqnarray}\nHere $f_M = M_{\\rm opt}^3 \\sin^3i\/(M_{\\rm NS}+M_{\\rm opt})^2$ is the\nmass function of the binary, and $f_0$ and $\\sigma_f$ are the measured\nvalue of $f_M$ and its uncertainty (obtained from X-ray timing\nobservations). Similarly, $K_0$, $\\sigma_K$ are the measured value of\n$K_{\\rm opt}$ and its uncertainty, $v_0$, $\\sigma_v$ are the measured\nvalue of $v_{\\rm rot}\\sin i$ and its uncertainty, and $\\theta_0$,\n$\\sigma_\\theta$ are the measured semi-duration of the eclipse\n$\\theta_e$ and its uncertainty. All these measurements are listed in\nRawls et al.\\ (2011) for the six systems of interest. Note that there\nis no measurement of $v_{\\rm rot} \\sin i$ for Her X-1, which leads to\nlarger uncertainties in the mass determination of this source.\n\nThe quantity $\\beta$ is equal to $R_{\\rm opt}\/R_L$, where $R_L$ is the\neffective radius of the Roche lobe of the secondary. For the\nintegration limits in equation (\\ref{eclipse}), we choose $\\beta_{\\rm\nmin}=0$, $\\beta_{\\rm max}=1$, $\\Omega_{\\rm min}=0$, $\\Omega_{\\rm\nmax}=2\\Omega_b$, where $\\Omega_b=2\\pi\/P_b$ is the mean angular\nvelocity of the binary orbit, and a sufficiently generous range of\n$M_{\\rm opt}$. As indicated in equation (\\ref{eclipse}), we assume a\nflat prior for each of the variables, though we have confirmed that\nthe results are not sensitive to this assumption. We calculate $P({\\rm\ndata}|M_{\\rm NS})$ for each of the six X-ray pulsars by computing the\nintegrals via a Monte Carlo method. We show in\nFigure~\\ref{fig:psr_slow} the resulting likelihoods of mass for each\nof the six systems and summarize the results in Table~6.\n\nWe consider an eccentric orbit in the case of Vela X-1 using the value\nof $e$ given in Rawls et al. (2011), and we assume a circular orbit\nfor the other five systems. Rawls et al. (2011) considered both\neccentric and circular orbits for 4U1538--52. However, the evidence\nfor eccentricity is not very strong. Moreover, the estimate of $K_{\\rm\nopt} =14.1\\pm1.1 ~{\\rm km\\,s^{-1}}$ that they obtain for their\neccentric orbit solution appears to be anomalously low (it does not\nfit the measured velocities very well --- see the lower left panel of\ntheir Fig. 9). For this reason, we consider only a circular orbit for\n4U1538--52, and we take $K_{\\rm opt} = 21.8\\pm3.8 ~{\\rm km\\,s^{-1}}$\nas estimated by Rawls et al. (2011). (The fit to the velocity data\nappears to be better in this case --- see the bottom right panel of\ntheir Fig. 9.)\n\n\\input{table6.tex}\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f5.ps}\n\\caption{The likelihoods $P_i({\\rm data}\\vert M_{\\rm NS})$ for \nthe eclipsing X-ray pulsars in high-mass X-ray binaries, which belong\nto category IV discussed in the text. This figure also includes the\nlikelihoods for the slow pulsars PSR~J1141$-$6545 and PSR~B2303+46,\nwhich belong to categories Ib and IIb, respectively, but have not been\nrecycled.}\n\\mbox{}\n\\label{fig:psr_slow} \n\\end{figure}\n\n{\\it (V) Accreting Bursting Neutron Stars.---\\\/} Spectroscopic\nobservations of a number of neutron stars during thermonuclear X-ray\nbursts have led to measurements of their masses and radii. This\ncategory includes 4U~1745$-$248 (\\\"Ozel et al.\\ 2009), 4U~1608$-$52\n(G\\\"uver et al.\\ 2010a), 4U 1820$-$30 (G\\\"uver et al.\\ 2010b), and\nKS~1731$-$260 (\\\"Ozel et al.\\ 2012). We convert the posterior\nlikelihood of mass and radius for each neutron star reported in\nthese analyses into a likelihood of mass by integrating over radius\nas\n\\begin{equation}\nP_i({\\rm data} \\vert M_{\\rm NS}) = C_i \\int dR_{\\rm NS} P_i({\\rm data} \n\\vert M_{\\rm NS}, R_{\\rm NS}).\n\\end{equation}\nThe resulting likelihoods are shown in Figure~\\ref{fig:bursters}. \n\nThe mass of a fifth burster, Cyg X-2, has been measured in a different\nmanner using the optical observations of its companion star (Orosz \\&\nKuulkers 1999). These observations yield the mass function of the\nbinary $f = 0.69 \\pm 0.03~M_\\odot$, a range of allowed mass ratios $0.3 < q <\n0.38$, as well as a measurement of the binary inclination $i= 61^\\circ\n\\pm 12^\\circ$. Using this information, we calculate the posterior\nlikelihood of the mass of the neutron star according to (\\\"Ozel et\nal.\\ 2010a)\n\\begin{equation}\nP_i({\\rm data} \\vert M_{\\rm NS}) = C_i \\int_{q_{\\rm min}}^{q_{\\rm max}} dq \n\\int_{(\\cos i)_{\\rm min}}^1 \\frac{d(\\cos i)}{1-(\\cos i)_{\\rm min}}\n\\exp \\left\\{- \\frac{[f_{0,i}-M_{\\rm NS} \\sin^3 i\/(1+q)^2]^2} {2 \\sigma_{f,i}^2}\n- \\frac{(i-i_0)^2}{2 \\sigma_i^2}\n\\right\\}\n\\label{eq:cygx2}\n\\end{equation}\nand plot it in Figure~\\ref{fig:bursters}.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f6.ps}\n\\caption{The likelihoods $P_i({\\rm data}\\vert M_{\\rm NS})$ for the accreting \nbursting neutron stars discussed as category V in the text. The likelihood \nfor Cyg~X-2 was inferred using optical observations of its companion star. }\n\\label{fig:bursters} \n\\mbox{}\n\\end{figure}\n\n\\section{The Intrinsic Distribution of Neutron Star Masses}\n\nThe range of neutron star masses that can be produced in\nastrophysically plausible scenarios is rather narrow, from $\\sim\n1.0~M_\\odot$ to $\\sim 2.5~M_\\odot$. The posterior likelihood functions\nfor the neutron star mass measurements we report in the previous\nsection are indeed within this range. Moreover, earlier studies of the\nneutron star masses indicated a narrowly peaked distribution (Finn\n1994; Thorsett \\& Chakrabarty 1999). For these reasons, we will model\nthe distribution of neutron star masses with a mean $M_0$ and a\ndispersion $\\sigma$, i.e.,\n\\begin{equation}\nP(M_{\\rm NS}; M_0,\\sigma) = \\frac{1}{\\sqrt{2 \\pi \\sigma^2}} \n\\exp\\left[-\\frac{(M_{\\rm NS}-M_0)^2}{2 \\sigma^2}\\right].\n\\label{eq:gauss}\n\\end{equation} \nOur goal in this section is to obtain the most likely values for the\nparameters $M_0$ and $\\sigma$ of this distribution that are consistent\nwith the measurements. \n\nUsing Bayes' theorem, we can write the probability $P(M_0,\\sigma \\vert\n{\\rm data})$ that measures the likelihood of the parameters of the\nneutron star mass distribution as\n\\begin{equation}\nP(M_0, \\sigma \\vert {\\rm data}) = C P({\\rm data} \\vert M_0,\\sigma) P(M_0) P(\\sigma),\n\\label{eq:bayes}\n\\end{equation}\nwhere $C$ is an appropriate normalization constant and $P(M_0)$ and\n$P(\\sigma)$ are the priors over the parameters of the mass\ndistribution. Hereafter, we will assume a flat prior over $M_0$\nbetween 1~$M_\\odot$ and $3~M_\\odot$ as well as a flat distribution\nover $\\sigma$ between zero and $1~M_\\odot$. \nWe repeated the analysis using a logarithmic prior distribution, which had\nindistiguishable effects on the final results. This is expected, given that\nthe inferred mass distributions are all very narrow and, therefore, are not\nsensitive to any weak priors.\n\nIn equation~(\\ref{eq:bayes}), the quantity $P({\\rm data} \\vert\nM_0,\\sigma)$ measures the posterior probability of having made a\nparticular set of observations for the ensemble of neutron stars given\nthe values of the parameters of the mass distribution. Under the \nassumption that each measurement is independent of all the others, \nwe calculate this quantity using \n\\begin{equation}\nP({\\rm data} \\vert M_0,\\sigma) = \\prod_i \\int dM_{\\rm NS} \nP_i({\\rm data} \\vert M_{\\rm NS}) P(M_{\\rm NS}; M_0, \\sigma).\n\\end{equation}\nThe only case where this assumption is not satisfied is for double\nneutron stars in category IIa, for which the mass measurements of the\npulsar and the companion neutron star are not independent. For these \nthree binary systems, we write instead \n\\begin{eqnarray}\nP({\\rm data} \\vert M_0,\\sigma) &=& \\prod_i^3 \\int dM_{\\rm tot} \n\\exp\\left[-\\frac{(M_{\\rm tot}-M_{\\rm tot,0})^2}{2 \n\\sigma_{M_{\\rm tot}}^2} \\right] \\nonumber \\\\\n&\\times&\n\\int dM_{\\rm NS} \nP_i({\\rm data} \\vert M_{\\rm NS}) P(M_{\\rm NS}; M_0, \\sigma) \nP(M_{\\rm tot}-M_{\\rm NS}; M_0, \\sigma).\n\\end{eqnarray}\n\nIn Section~2, we divided the neutron star mass measurements not only\naccording to the technique by which these measurements were obtained\nbut also by the type of companion and the spin period of the neutron\nstar. We carried out the latter division in anticipation of the fact\nthat fast and slow pulsars are drawn from different parent\npopulations, i.e., from those which have and have not experienced\nsignificant mass accretion phases. In the following, we derive the\nparameters of the intrinsic mass distributions for each of those\npopulations separately. We then address the extent to which the\nparticular evolutionary channels followed by each type of neutron star\nleaves a measurable imprint on the mass distribution.\n\n\\subsection{Neutron Stars at or Near their Birth Masses}\n\nThe low spin periods of a number of pulsars in our sample are\nindicative of mild or even no recycling due to mass accretion. We,\ntherefore, consider the masses of neutron stars in this population\nlikely to be very near their birth values. This sample includes\ncategories Ia and IIa for double neutron stars, category IV for\naccreting pulsars with primarily high-mass companions, as well as one\npulsar each in categories Ib (PSR~J1141$-$6545) and IIb\n(PSR~B2303+46).\n\nWe first study the underlying mass distribution of the double neutron\nstars. The mass measurements in these systems have by far the smallest\nerrors, which can dominate the parameter estimation of the mass\ndistribution of the total ensemble. Furthermore, these binaries have\nfollowed a very particular and highly selective evolutionary path,\nwhich may be evident in their mass distribution. We group the\nremaining sources together as a second sample that consists of neutron\nstars likely to be near their birth masses.\n\n\\subsubsection{Double Neutron Stars}\n\nFigure~\\ref{fig:prob_dns} shows the 68\\% and 95\\% confidence contours\nover the parameters of the intrinsic Gaussian distribution that is\nconsistent with the observed masses of nine double neutron stars. The\nmost likely value of the mean of the Gaussian distribution is\n1.33~$M_\\odot$ and that of the dispersion is 0.05~$M_\\odot$. It is\nevident from the figure that the uncertainties in the parameters of\nthe underlying distribution are very small: the 68\\% errors are\n0.03~$M_\\odot$ in both parameters.\n\nOur results are in agreement with the distribution reported by\nThorsett \\& Chakrabarty (1999), who found a mean of 1.35~$M_\\odot$ and\na dispersion of 0.04~$M_\\odot$, and with the more recent results of\nKiziltan et al.\\ (2010) for this category.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f7.ps}\n\\caption{The confidence contours over the parameters of a Gaussian \ndistribution for the double neutron stars.}\n\\mbox{}\n\\label{fig:prob_dns} \n\\end{figure}\n\nAn interesting question we can address with the sample of double\nneutron stars is whether the two members of each binary system are\ndrawn from the same underlying population. To study this, we divided\nthe sample into one of pulsars and one of the companions. Note that\nfor the double pulsar, we assigned the faster pulsar to the ``pulsar''\nand the slower to the ``companion'' categories. Repeating the above\ninference for these two subgroups individually, we found that the most\nlikely parameters of the mass distribution for the pulsars are $M_0 =\n1.35~M_\\odot$ and $\\sigma = 0.05~M_\\odot$, whereas for the companions\nthey are $M_0 = 1.32~M_\\odot$ and $\\sigma = 0.05~M_\\odot$. These\nparameters are the same within the 68\\% confidence ranges of each.\n\nOne further question about the double neutron star population concerns\nthe mass ratio $q$ in each binary. The mass ratios of the neutron\nstars in two of the six binaries with well determined masses are\nwithin one part in $\\sim 5 \\times 10^{-3}$ of unity. This is an order\nof magnitude smaller than the most likely dispersion of the underlying\nmass distribution. We explored whether the distribution of observed\nmass ratios is consistent with the pulsar and the companion being\ndrawn independently from a Gaussian distribution with the parameters\nwe determined above.\n\nThe posterior likelihood of observing a binary with pulsar and\ncompanion masses of $M_{\\rm psr}$ and $M_{\\rm c}$, respectively, is\ngiven by\n\\begin{equation}\nP(M_{\\rm psr}, M_{\\rm c}) dM_{\\rm psr} dM_{\\rm c} = C \\exp \n\\left[-\\frac{(M_{\\rm psr}-M_0)^2}{2\\sigma^2}-\\frac{(M_{\\rm c}-M_0)^2}{2\\sigma^2}\\right]\ndM_{\\rm psr} dM_{\\rm c}\n\\end{equation}\nTo convert this into a distribution over the mass ratio, we set \n$q \\equiv \\min(M_{\\rm psr}\/M_{\\rm c}, M_{\\rm c}\/M_{\\rm psr})$ and write \n\\begin{eqnarray}\nP(q)dq &=& \\int_{M_{\\rm psr}} P(M_{\\rm psr},q) dM_{\\rm psr} dq = \n\\int_{M_{\\rm psr}} P(M_{\\rm psr},M_{\\rm c}) \n\\frac{dM_{\\rm c}}{dq} dM_{\\rm psr} dq \\nonumber \\\\\n &=& \\int_{M_{\\rm psr}} \\frac{2\\;C}{M_{\\rm psr}} \n\\exp \\left[-\\frac{(M_{\\rm psr}-M_0)^2}{2\\sigma^2}-\n\\frac{(q\\;M_{\\rm psr}-M_0)^2}{2\\sigma^2}\\right]\ndM_{\\rm psr} dq\\;.\n\\end{eqnarray}\nIn Figure~\\ref{fig:dns_massratio}, we compare the cumulative\nlikelihood of the mass ratio\n\\begin{equation}\nC(q>q_0)=\\int_{q_0}^1 P(q)dq\\;,\n\\end{equation}\ncalculated for the most likely values of the parameters of the\nGaussian, to the cumulative distribution of the observed mass ratios\nof the 6 double neutron stars with well determined masses. The\nsimilarity between the two distributions is striking and demonstrates\nthat the pulsar and the companion in each of the double neutron stars\nare consistent with having been drawn independently from the same\nnarrow distribution of masses.\n\nWe also explored whether the observed distribution of mass ratios is\nconsistent with the predicted cumulative distribution for neutron star\npairs drawn independently from the double Gaussian distribution\nsuggested by Schwab et al.\\ (2010; their eq.~(1)). The result is shown\nas a green line in Figure~\\ref{fig:dns_massratio}. The width of the\nindividual components in the bimodal distribution is significantly\nnarrower than the width of the single Gaussian that we infer here.\nThis leads to a larger fraction of double neutron stars with mass\nratios closer to unity for the bimodal distribution, which is not in\nagreement with the observed sample.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f8.ps}\n\\caption{The histogram shows the cumulative mass ratio distribution for \nthe six double neutron stars with precise mass measurements. The red\nline shows the predicted cumulative distribution for neutron star\npairs drawn independently from a single Gaussian distribution with a\ncentral value and a dispersion equal to the most likely parameters\nshown in Figure~\\ref{fig:prob_dns}. The green line shows the predicted\ncumulative distribution for neutron star pairs drawn independently\nfrom the double Gaussian distribution suggested by Schwab et al.\\\n(2010). The observed distribution of mass ratios is in agreement with\na mass distribution represented by a single Gaussian. Note that, for\nconsistency, we show in this figure the mass ratio histogram generated\nfrom the data used by Schwab et al.\\ (2010).}\n\\mbox{}\n\\label{fig:dns_massratio} \n\\end{figure}\n\n\\subsubsection{Accreting and Slow Pulsars}\n\nThe second subgroup consists of neutron stars accreting from high-mass\ncompanions and slow pulsars, which are likely to be near their birth\nmasses. To infer the neutron star mass distribution, we will use both\nthe numerical results of Rawls et al.\\ (2011) as well as our analytic\nresults discussed in \\S 2, in which the information regarding\nellipsoidal variations in the lightcurves was not taken into account.\n\nIn Figure~\\ref{fig:prob_slow_num}, we show the 68\\% and 95\\%\nconfidence limits on the Gaussian parameters of the underlying mass\ndistribution for the accreting and slow pulsars, using the numerical\nresults of Rawls et al.\\ (2011). For comparison, we overplot the\nequivalent confidence contours for the double neutron stars. There is\na small but statistically insignificant shift in the central mass of\nthe Gaussian between the two populations. On the other hand, the\nGaussian dispersions between the two populations are different to a\nhigh statistical significance. In other words, even though both of\nthese populations are believed to represent neutron stars near their\nbirth masses, the double neutron stars are drawn from a significantly\nnarrower distribution of masses. The most likely values of the central\nmass and dispersion for the accreting and slow pulsars are\n1.28~$M_\\odot$ and 0.24~$M_\\odot$, respectively.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f9.ps}\n\\caption{The confidence contours over the parameters of a Gaussian \ndistribution for the accreting and slow pulsars using the numerical\ndata from Rawls et al.\\ (2011). The confidence contours for the double\nneutron stars are also shown for comparison. Even though both\npopulations are thought to have masses near their birth masses, the\ndispersion of double neutron star masses is significantly smaller than\nthat of the accreting and slow pulsars.}\n\\mbox{}\n\\label{fig:prob_slow_num} \n\\end{figure}\n\nThe Rawls et al.\\ (2011) analysis depends on a synthesis of a large\nnumber of spectroscopic and photometric measurements of the binaries\nthat are used to infer the binary parameters. Fitting these\nobservations, and especially taking into account the ellipsoidal\nvariations requires complex models of the shape of the companion star\nand of the relative contribution of light from the accretion\ndisk. Moreover, the photometric observations of the ellipsoidal\nvariations typically are the lowest signal-to-noise components of the\nmass measurements. In order to assess the possible influence of the\nmodeling of the ellipsoidal variations on our results, we also infer\nthe underlying mass distribution using our analytical posterior\nprobabilities discussed in \\S 3. Figure~\\ref{fig:prob_slow_anlytc}\nshows the resulting 68\\% and the 95\\% confidence contours for the\nparameters of the mass distribution. In this case, the most likely\nvalues of the central mass and dispersion for the accreting and slow\npulsars are 1.24~$M_\\odot$ and 0.20~$M_\\odot$, respectively. Comparing\nthese to the distribution inferred from the numerical results of Rawls\net al.\\ (2011), we see that the central mass remains unchanged, but\nthe dispersion becomes less constrained and is even statistically\nconsistent with that of double neutron stars.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f10.ps}\n\\caption{The confidence contours over the parameters of a Gaussian \ndistribution for the accreting and slow pulsars using the analytic\nmass measurements discussed in the text. The confidence contours for\nthe double neutron stars are also shown for comparison. }\n\\mbox{}\n\\label{fig:prob_slow_anlytc} \n\\end{figure}\n\n\\subsection{Recycled Neutron Stars}\n\nWe now focus on the subgroup of neutron stars which have been recycled\nthrough extended mass accretion. Neutron stars with white dwarf\ncompanions, millisecond pulsars, as well as neutron stars in low-mass\nX-ray binaries currently undergoing accretion (categories Ib, IIb,\nIII, and V) belong to this group. Figure~\\ref{fig:prob_recycl} shows\nthe confidence contours over the parameters of the Gaussian\ndistribution for the recycled neutron stars. The most likely value of\nthe central mass is 1.48~$M_\\odot$ and of the dispersion is\n0.20~$M_\\odot$. Both the dispersion and the mean are similar to those\nfound by Kiziltan et al.\\ (2010) within statistical uncertainties. The\nuncertainties in the mean value quoted by Kiziltan et al.\\ (2010),\nhowever, are significantly smaller than those shown in\nFigure~\\ref{fig:prob_recycl} (c.f. Figure~3 of Kiziltan et al.\\ 2010).\n\nThere are two main differences between our study and that of Kiziltan\net al. (2010) with regard to the recycled neutron star sample. First,\nwe make use of the detailed posterior likelihood for each mass\nmeasurement, whereas Kiziltan et al.\\ (2010) appear to have\napproximated them with asymmetric Gaussians. Second, we include in our\nsample neutron stars in low-mass X-ray binaries for which mass\nmeasurements were performed mostly spectroscopically and typically\nhave larger uncertainties.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f11.ps}\n\\caption{The confidence contours over the parameters of a Gaussian \ndistribution for the recycled neutron stars. The confidence contours\nfor the double neutron stars are also shown for comparison. As\nexpected, the recycled neutron stars have on average larger masses\nthan those in double neutron stars.}\n\\mbox{}\n\\label{fig:prob_recycl} \n\\end{figure}\n\nIn order to assess the sensitivity of our results to the inclusion of\nthe accreting neutron stars, we repeat the inference of the mass\ndistribution parameters using only the radio pulsars in categories Ib\nand IIb. In Figure~\\ref{fig:prob_fastpsr} we show the resulting\nconfidence contours. The difference with the entire sample is minimal:\nthe most likely mean value and the dispersion are 1.46~$M_\\odot$ and\n0.21~$M_\\odot$, respectively. We, therefore, attribute the small\ndifference with the Kiziltan et al.\\ (2010) results to our handling of\nthe posterior likelihood distributions for each measurement.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f12.ps}\n\\caption{The confidence contours over the parameters of a Gaussian \ndistribution for a subgroup of the recycled neutron stars that\nincludes only pulsars in orbit around white dwarfs. Considering only\nthese sources with dynamical mass measurements does not alter the\nresults shown in Figure~\\ref{fig:prob_recycl}. }\n\\mbox{}\n\\label{fig:prob_fastpsr} \n\\end{figure}\n\n\n\\section{Discussion}\n\nIn this paper, we investigated the distribution of neutron star masses\nin different types of binary systems and at different stages of\nevolution based on currently available measurements. We summarize the\nneutron star mass measurements and their uncertainties in each\nsubgroup in Figure~\\ref{fig:nsmasses} and compare them to those of\nblack holes in Figure~\\ref{fig:bhmasses} (compiled and analyzed in\n\\\"Ozel et al.\\ 2010a). In these figures, the error bars correspond \nto a 68\\% confidence level calculated from the detailed likelihood\ndistribution presented for each subgroup of sources in \\S 2.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f13.ps}\n\\caption{The masses of neutron stars measured in double neutron stars \n(magenta; categories Ia and IIa), in eclipsing binaries with primarily\nhigh mass companions (cyan; category IV; these are the numerical\nvalues from Rawls et al.\\ 2011 given in column~2 of Table~6), with\nwhite dwarf companions (gold; categories Ib and IIb), with optical\nobservations of the white dwarf companions (green; category III), and\nin accreting bursters (purple; category V). }\n\\mbox{}\n\\label{fig:nsmasses} \n\\end{figure}\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f14.ps}\n\\caption{The measured masses of Galactic black holes \n(after \\\"Ozel et al. 2010a).}\n\\label{fig:bhmasses} \n\\mbox{}\n\\end{figure}\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.6]{f15a.ps}\n \\includegraphics[scale=0.6]{f15b.ps}\n\\caption{The inferred mass distributions for the different populations \nof neutron stars (top) and black holes (bottom) discussed in the text.\nThe dashed lines correspond to the most likely values of the\nparameters. For the different neutron star populations these are:\n$M_0=1.33 M_\\odot $ and $\\sigma=0.05 M_\\odot$ for the double neutron\nstars, $M_0=1.28 M_\\odot $ and $\\sigma=0.24 M_\\odot$ for the other\nneutron stars near their birth masses, and $M_0=1.48 M_\\odot $ and\n$\\sigma=0.20 M_\\odot$ for the recycled neutron stars. For the case of\nblack holes, we used the exponential distribution with a low mass\ncut-off at $M_{\\rm c}=6.32 M_\\odot$ and a scale of $M_{\\rm scale}=1.61\nM_\\odot$ obtained in \\\"Ozel et al.\\ (2010a). The solid lines represent\nthe weighted mass distributions for each population, for which\nappropriate fitting formulae are given in the Appendix. The\ndistributions for the case of black holes have been scaled up by a\nfactor of three for clarity. }\n\\mbox{}\n\\label{fig:weighted} \n\\end{figure}\n\nIn the top panel of Figure~\\ref{fig:weighted}, we show the inferred\nmass distributions of the various neutron star populations discussed\nin Section~2. For each population, we present two different\ndistributions. The dashed lines correspond to the most likely\nparameters of the underlying distributions inferred in Section~3. Each\nsolid line represents the weighted distribution over the central mass\nand dispersion for each population. We compute this weighted\ndistribution as\n\\begin{equation}\nP_w(M_{\\rm NS})=\\int dM_0 \\int d\\sigma P(M_{\\rm NS}; M_0,\\sigma)\nP(M_0,\\sigma \\vert {\\rm data}), \n\\end{equation}\nwhere $P(M_{\\rm NS}; M_0,\\sigma)$ and $P(M_0,\\sigma \\vert {\\rm data})$\nis given by equations~\\ref{eq:gauss} and \\ref{eq:bayes},\nrespectively. In the Appendix, we provide approximate analytic fitting\nformulae for these weighted distributions for each population.\n\nIn the bottom panel of Figure~\\ref{fig:weighted}, we compare the\ninferred mass distribution for recycled neutron stars to that of black\nholes reported in \\\"Ozel et al.\\ (2010a). For the latter, we use the\nexponential model with a lower mass cutoff given by\n\\begin{equation}\nP(M_{\\rm BH}; M_{\\rm scale}, M_{\\rm c}) = \n\\frac{\\exp(M_{\\rm c}\/M_{\\rm scale})}{M_{\\rm scale}}\n\\times\\left\\{\n\\begin{array}{ll}\n\\exp(-M_{\\rm BH}\/M_{\\rm scale}), & M_{\\rm BH}>M_{\\rm c} \\\\\n0, & M_{\\rm BH}\\leq M_{\\rm c}\n\\end{array}\n\\right..\n\\end{equation}\nThe most likely values for the parameters of this distribution are\n$M_{\\rm scale}=1.61 M_\\odot$ and $M_{\\rm c} = 6.32 M_\\odot$. In the\nsame panel, we also include the appropriate weighted distribution for\nthe black holes, where we carried out the integration over the\nposterior likelihood of the parameters $M_{\\rm scale}$ and $M_{\\rm\nc}$; we provide an analytic fitting formula for the weighted\ndistribution in the Appendix. This panel highlights the substantial\nmass gap that exists between the black hole population and even the\nheaviest neutron star population (see the discussion in\n\\\"Ozel et al.\\ 2010a and Farr et al.\\ 2011).\n\nWithin the neutron star population, it is evident from these figures\nthat the mass distribution of double neutron star systems is different\nthan those observed in other binary systems, which include both\nneutron stars near their birth masses as well as neutron stars that\nexperienced significant accretion episodes. Indeed, the most likely\nvalues of the mean mass and the dispersion we derived for these\npopulations using the Bayesian inference technique discussed in \\S 3\nare $1.33 \\pm 0.05~M_\\odot$ for double neutron stars, in contrast to\n$1.28 \\pm 0.24~M_\\odot$ for other neutron stars near their birth\nmasses and $1.48 \\pm 0.20~M_\\odot$ for recycled neutron stars. Note\nthat the uncertainties in both the mean mass and the dispersion for\nall of these subgroups are shown in Figures~9-12.\n\nThe narrowness of the mass distribution of double neutron stars is\ndifficult to account for within the current understanding of neutron\nstar formation mechanisms. One possible way to generate a narrow\ndistribution is via electron capture supernovae in ONeMg white dwarfs.\nThe onset of such a supernova occurs at a particular density\nthreshold, which corresponds to a pre-collapse mass of the white dwarf\nin the narrow range $1.36-1.38~M_\\odot$ for different temperatures and\ncompositions (Podsiadlowski et al.\\ 2005). Taking into account a\nbinding mass given by the approximate formula (Lattimer \\& Yahil 1989)\n\\begin{equation}\nE_B = 0.084 \\left(\\frac{M}{M_\\odot}\\right)^2 M_\\odot\n\\end{equation}\nthe gravitational masses of the outcomes of electron capture\nsupernovae become $1.2-1.22~M_\\odot$. This range of masses is compared\nto the parameters of the underlying distributions of double neutron\nstars as well as of the other neutron stars near their birth masses in\nFigure~\\ref{fig:nsmass_theory}. Even though the electron capture\nsupernovae are capable of producing a narrow range of neutron star\nmasses, the mean of the expected distribution is inconsistent with\nthat of double neutron stars to a high confidence level.\n\n\\begin{figure}\n\\centering\n \\includegraphics[scale=0.75]{f16.ps}\n\\caption{The distribution of neutron star masses at birth expected from \ntheoretical calculations, compared to the observed mass distribution\nof neutron stars that have not accreted significant amounts of mass\n(labeled accreting and slow pulsars) as well as to that of double\nneutron stars. The parameters of the distribution of the former\nsubgroup are consistent with expectations from core-collapse supernova\nand fallback, while the observed double neutron star distribution is \nsignificantly narrower than what is expected.} \n\\mbox{}\n\\label{fig:nsmass_theory} \n\\end{figure}\n\nIn the case of core collapse supernovae, the mean mass of the neutron\nstars produced in the absence of fallback is also expected to be\nsignificantly smaller than that inferred from the double neutron\nstars. We can estimate this mass assuming that the core of the\npre-supernova star collapses when it reaches its Chandrasekhar limit.\nConsidering an electron fraction of $Y_e = 0.42-0.48$, which is\nappropriate for the cores of presupernova stars (Timmes et al.\\ 1996), \nthe Chandrasekhar mass \n\\begin{equation}\nM_{\\rm Ch} = 5.83 Y^2_e\n\\end{equation}\nfalls in the range $1.15-1.34~M_\\odot$. Taking into account the\ngravitational binding energy, the expected range of birth masses for\nneutron stars from core-collapse supernovae is $1.06-1.22~M_\\odot$.\n\nFallback of matter during and immediately following the supernova\nexplosion can naturally lead to neutron stars more massive than the\ncores of the progenitor stars. At the same time, the stochastic nature\nof fallback necessarily leads to an increased dispersion of neutron\nstar masses. In Figure~\\ref{fig:nsmass_theory}, we show the evolution\nof the expected dispersion with central mass assuming that a fallback\nof baryonic mass $\\Delta M_f$ introduces a dispersion of the baryonic\nmass of the neutron star of the same magnitude. (Note that in\nFig.~\\ref{fig:nsmass_theory}, we plot the corresponding gravitational\nmass for the neutron stars.) This simple analytical estimate is in\nagreement with the detailed numerical calculations of Zhang et al.\\\n(2008), which are also shown in the figure. The green hatched region\noutlines the results for different compositions, explosion energies,\nand locations of the pistons for an assumed maximum neutron star mass\nof $2~M_\\odot$. Allowing sufficient fallback to account for the mean\nvalue of the double neutron star masses introduces a dispersion that\nis significantly larger than the observed one. In contrast, the\ninferred mean and dispersion of the mass distribution of other neutron\nstars believed to be near their birth masses (labeled accreting and\nslow pulsars) are in agreement with theoretical expectations of core\ncollapse supernovae.\n\nConsidering a bimodal mass distribution for the double neutron stars,\nas in Schwab et al.\\ (2010), aggravates these challenges. First, as we\nshowed in \\S 3, the cumulative likelihood of the mass ratio for such\na distribution does not agree with the cumulative distribution of the\nobserved mass ratios for the double neutron stars. Second, the\ndispersion in the two components becomes even smaller, $\\simeq\n0.025~M_\\odot$, making the higher mass component even less consistent\nwith the expectations of the core collapse supernovae. All of these\narguments lead to the conclusion that the mass distribution of double\nneutron stars is peculiar and perhaps related to the particular\nevolutionary history that leads to their formation.\n\nThe masses of the population of recycled neutron stars, which include\nfast pulsars with white dwarf companions as well as accreting\nbursters, are consistent with them having undergone extended periods\nof accretion. On average, recycled neutron stars are more massive by\n$\\approx 0.2 M_\\odot$ compared to other accreting and slow\npulsars. Such a mass increase is more than adequate to recycle these\npulsars to millisecond periods. Indeed, assuming that the mass is\ntransferred onto the neutron star via an accretion disk that is\nmagnetically truncated at the corotation radius\n\\begin{equation}\nR_c = \\left(\\frac{GM}{4 \\pi^2 \\nu^2_{\\rm s}}\\right)^{1\/3}, \n\\end{equation}\nwhere $\\nu_{\\rm s}$ is the spin frequency of the neutron star, \nthe angular momentum transferred per unit mass is \n\\begin{equation}\nl= (G M R_c)^{1\/2} = \\left(\\frac{G^2 M^2}{2 \\pi \\nu_{\\rm s}}\\right)^{1\/3}.\n\\end{equation}\nAfter accreting mass $\\Delta M$, the neutron star acquires an angular \nmomentum $\\Delta M \\cdot l$. Equating this to the spin angular momentum of the \nrecycled pulsar $L=2 \\pi I \\nu_{\\rm s}$, where $I$ is its moment of inertia, \nallows us to calculate the mass required to spin up the pulsar as \n\\begin{equation}\n\\Delta M = I (GM)^{-2\/3} (2 \\pi \\nu_{\\rm s})^{4\/3} \n= 0.034 \\left(\\frac{\\nu_{\\rm s}}{300~{\\rm Hz}}\\right)^{4\/3}\n\\left(\\frac{M}{1.48~M_\\odot}\\right)^{-2\/3}\n\\left(\\frac{I}{10^{45}~{\\rm g}~{\\rm cm}^2}\\right)\\;M_\\odot. \n\\end{equation}\n\nIt is interesting that the most likely value of the mean mass of the\nrecycled pulsars is significantly smaller than the $2~M_\\odot$ lower\nbound on the maximum mass of a neutron star (Demorest et al.\\ 2010;\n\\\"Ozel et al.\\ 2010b) as well as the average mass of recycled\nneutron stars predicted by population synthesis studies (e.g., Pfahl,\nRappaport, \\& Podsiadlowski 2003; Lin et al.\\ 2011). This conclusion\ncan be used to refine models of low-mass X-ray binary evolution.\nFurthermore, our analysis shows that a very small fraction of neutron\nstars reach masses comparable to the maximum possible neutron star\nmass and collapse into black holes. Therefore, this channel does not\ncontribute significantly to a putative but still undetected population\nof low-mass black holes in the Galaxy (see the discussion in \\\"Ozel et\nal.\\ 2010a).\n\n\\acknowledgments\n\nF.\\\"O., D.P., and R.N. thank the Institute of Astronomy, where a large\nfraction of this work was carried out, for their hospitality. We thank\nScott Ransom for useful discussions as well as Ingrid Stairs, Michael\nKramer, and Matthew Bailes for detailed comments on the manuscript. We\nthank Laura Kasian for sharing unpublished data. We gratefully\nacknowledge support from NSF grant AST-1108753, NASA ADAP grant\nNNX10AE89G, NSF CAREER award NSF 0746549, and Chandra Theory grant\nTMO-11003X for this work.\n\n\n\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\\label{sec:introduction}\n\nHydrogenic impurities are a central aspect of semiconductor physics and technology.\\cite{pan,yu-car}\nIn recent years, impurity states have been proposed as potential qubits in quantum information \ndevices.\\cite{kan,koi-hu-das}\nIn binary III-V and II-VI semiconductors, the spin-orbit interaction can play an important role\nin the electronic structure of confined electronic states.\\cite{win,int-tam-wei}\nAt the level of the effective-mass approximation, a hydrogenic donor impurity and the hydrogen \natom are almost completely analogous quantum systems.\\cite{koh}\nThus, it is in principle a simple matter to describe at that level the effect of the so-called \nextrinsic or Rashba spin-orbit interaction in the impurity states.\nThis contribution to the spin-orbit interaction in local external potentials in solids is, \nin this context, analogous to the spin-orbit coupling in the hydrogen atom.\nHowever, care must be taken due to the necessary renormalization of the spin-orbit coupling \nconstant and to possible modifications to the spin-orbit formula due to anisotropies of the \ncrystal structure.\nOn the other hand, an additional contribution to the spin-orbit interaction, which is specific \nto the solid-state context must be taken into consideration.\nThis so-called intrinsic or Dresselhaus contribution\\cite{dre} is present in the bulk of\nthe material and reappears in the envelope-function approximation theory of external, \nmesoscopic, potentials, like the one caused by a ionized donor. \n\n\nIn this article we calculate the energy levels of the $n=2$ shell of hydrogenic impurities\nof semiconductors with wurtzite crystal structure in the presence of these two contributions \nto the spin-orbit interaction.\nWhile, as mentioned above, the calculation of the extrinsic contribution is in a sense \na straightforward application of the well-known formulas for the hydrogen atom, we need to\nfill a theoretical gap, caused by the anisotropy of the wurtzite crystal structure.\nTo that effect, we obtain here an expression for the extrinsic spin-orbit interaction\nin wurtzite semiconductors and the effective spin-orbit coupling constant for conduction\nband electrons.\nThe calculation of the effects of the intrinsic spin-orbit interaction in this context\nis entirely new and it is considered acting alone and in combination with the extrinsic\ncontribution.\n\n\nThe article is organized as follows.\nIn Section II we introduce the system and study the effect of the intrinsic spin-orbit interaction.\nIn Section III we derive the effective Hamiltonian of the extrinsic spin-orbit interaction using \nthe Foldy-Wouthuysen transformation adapted to this context \nand in Section IV we calculate the energy corrections that it produces.\nIn Section V we study the combined effect of both spin-orbit interactions, and in Section VI we\nprovide the concluding remarks.\n\n\n\n\\section{Intrinsic spin-orbit interaction in the hydrogenic impurity}\n\\label{sec:efa-impurity}\n\nWe consider an electron bound to a hydrogenic donor impurity in a bulk semiconductor with\nwurtzite crystal structure.\nWorking at the level of the envelope-function approximation (EFA), both the intrinsic \nand the extrinsic spin-orbit couplings appear in the Hamiltonian:\n\\begin{equation} \n H = H_{\\text{0}} + H_{\\text{int}} + H_{\\text{ext}},\n\\label{eq:Htot}\n\\end{equation}\nwhere\n\\begin{equation} \n H_0 = \\frac{p^2}{2 m^*} + V(\\mathbf{r}).\n\\label{eq:Hzero} \n\\end{equation}\nHere $V(\\mathbf{r}) = - e^2\/\\epsilon r$ is the effective Coulomb potential of the electron\nbound to the ionized donor impurity.\nWe assume the effective mass $m^*$ and the dielectric constant $\\epsilon$ to be isotropic, \nthereby preserving the spherical symmetry of the hydrogenic Hamiltonian $H_0$.\n\n\nThe extrinsic spin-orbit coupling, $H_{\\text{ext}}$, will be discussed below, and for the\nmoment we focus on the intrinsic contribution, $H_{\\text{int}}$, which for semiconductors\nwith wurtzite crystal structure is given by \\cite{voo-wil-car, fu-wu,wan-wu-tsa}\n\\begin{equation}\n H_{\\text{int}} = \\alpha \\,( \\sigma_x k_y - \\, \\sigma_y k_x) \n + \\gamma \\,(b \\,k_z^{2}-\\, k_{\\parallel}^2)\n ( \\sigma_x k_y - \\, \\sigma_y k_x),\n\\label{eq:int_completa}\n\\end{equation}\nwhere $k_{\\parallel}^2=k_x^2+k_y^2$ and $\\alpha$, $b$, and $ \\gamma$ are material-dependent \nparameters which are obtained experimentally or via ab-initio calculations.\n$\\boldsymbol{\\hat{\\sigma}}=(\\sigma_x,\\sigma_y,\\sigma_z)$ are the Pauli matrices.\nWhile $\\alpha$ and $\\gamma$ can vary considerably between different materials,\n$b$ is roughly universal and close to 4 for all materials.\nNote that $H_{\\text{int}}$ has two parts, one of them linear and the other one cubic \nin the wavevector $k$.\nThe cubic-in-k term displays an anisotropy between the $z$-direction and the directions \nin the $xy$-plane.\nThis anisotropy and the presence of the linear term distinguish the intrinsic spin-orbit\nHamiltonian of wurtzite semiconductors from the Dresselhaus coupling of zincblende \nsemiconductors.\\cite{dre}\n\n\nThe hydrogenic Hamiltonian $H_0$ has the renormalized eigenvalues $E_n^*=-E_R^*\/n^2$,\nwhere $E_R^*=m^{\\ast} e^4\/2\\epsilon^2\\hbar^2$ is the effective Rydberg energy.\nThe aim of this study is to obtain the corrections to the $E_2$ energy level due to\nthe intrinsic and extrinsic spin-orbit Hamiltonians.\nWe will work at the level of first-order perturbation theory, which is adequate due\nto the smallness of the spin-orbit couplings compared to the separation of the \nbare $E_n^*$ levels.\nIn order to diagonalize the intrinsic Hamiltonian $H_{\\text{int}}$ in the $E_2$ subspace\nwe use the basis of hydrogenic eigenstates of $\\{L^2,L_z,S^2,S_z\\}$, given by\n\\begin{eqnarray}\n && \\psi_{200 \\eta} = \\left( \\frac{1}{32\\pi a^{*3}} \\right)^{\\frac{1}{2}}\n \\left(2 - \\frac{r}{a^*} \\right) \n e^{-r\/2a^*} |\\eta\\rangle \\nonumber \\\\\n && \\psi_{211 \\eta} = \\left( \\frac{1}{64\\pi a^{*3}} \\right)^{\\frac{1}{2}} \n \\frac{(x + i y)}{a^*} \\,\n e^{-r\/2a^*} |\\eta\\rangle \\nonumber \\\\\n && \\psi_{210 \\eta} = \\left( \\frac{1}{32\\pi a^{*3}} \\right)^{\\frac{1}{2}} \\frac{z}{a^*} \\,\n e^{-r\/2a^*} |\\eta\\rangle \\nonumber \\\\\n && \\psi_{21-1 \\eta} = -\\left( \\frac{1}{64\\pi a^{*3}} \\right)^{\\frac{1}{2}} \n \\frac{(x - i y)}{a^*} \\,\n e^{-r\/2a^*} |\\eta\\rangle ,\n\\end{eqnarray}\nwhere $\\eta= \\{\\uparrow, \\downarrow\\}$.\nThe matrix elements of the linear-in-$k$ terms of $H_{\\text{int}}$ in this basis are zero;\nonly the cubic-in-$k$ terms contribute.\nOrdering the basis states as: \n$|200\\uparrow\\rangle$, $|200\\downarrow\\rangle$, \n$|211\\uparrow\\rangle$, $|211\\downarrow\\rangle$, $|210\\uparrow\\rangle$, \n$|210\\downarrow\\rangle$, $|21,-1\\uparrow\\rangle$, and $|21,-1\\downarrow\\rangle$, \nthe matrix of $H_{\\text{int}}$ in the $n=2$ subspace is\n\\begin{equation}\n \\bar{\\bar H}_{\\text{int}} \\,= \\left(\\begin{matrix} \n \n 0 & 0 & 0 & A & 0 & 0 & 0 & 0 \\\\\n \n 0 & 0 & 0 & 0 & 0 & 0 & B & 0 \\\\\n \n 0 & 0 & 0 & 0 & 0 & 0 & 0 & 0 \\\\\n \n A & 0 & 0 & 0 & 0 & 0 & 0 & 0\\\\\n \n 0 & 0 & 0 & 0 & 0 & 0 & 0 & 0\\\\\n \n 0 & 0 & 0 & 0 & 0 & 0 & 0 & 0 \\\\\n \n 0 & B & 0 & 0 & 0 & 0 & 0 & 0 \\\\\n \n 0 & 0 & 0 & 0 & 0 & 0 & 0 & 0 \\\\\n \\end{matrix}\\right) ,\n\\end{equation}\nwhere\n\\begin{eqnarray}\n && A = \\frac{\\gamma}{32\\sqrt2 \\, a^{*3}} \\left(\\frac{14}{15} b + \\frac{133}{60} \\right) \\nonumber \\\\\n && B = \\frac{\\gamma}{32\\sqrt2 \\, a^{*3}} \\left(\\frac{62}{15} b + \\frac{433}{60} \\right) .\n \\label{eq:AB}\n\\end{eqnarray} \nThe secular equation for $\\bar{\\bar H}_{\\text{int}}$, \n$\\text{det} \\left(\\bar{\\bar H}_{\\text{int}} - \\varepsilon \\bar{\\bar I} \\right) = 0 $, \nyields\n\\begin{equation}\n \\varepsilon^{4} - \\varepsilon^{2}(A^2 + B^2) + (A\\,\\, B)^2 = 0 .\n \\label{eq:intAB}\n\\end{equation}\nThe eigenvalues are then $\\varepsilon_{1,2} = \\pm A $ and $\\varepsilon_{3,4} = \\pm B $.\nThe other four eigenvalues are degenerate and equal to zero.\nIn Table I we present the non-zero energy corrections for the materials GaN, ZnO, InN and AlN, \nalong with their $\\gamma$ and $b$ parameters.\nIn the last two columns we present the energy splittings $2A$ and $2B$ as percentages\nof the unperturbed energy $E_2^*$.\n\n\n\\begin{table}[ht]\n\\label{tab:int}\n\\begin{center}\n\\begin{tabular} {| l | c | c | c | c | c| c | c |}\n\\hline\n & $\\gamma $ & $ b $ & $ \\varepsilon_{1,2} $ & $ \\varepsilon_{3,4} $ & $ E_2^{*} $ & $2A\/E_2^{\\ast} $ & $2B\/E_2^{\\ast} $ \\\\\n &[$\\text{meV} \\text{\\AA}^3$] & & [ $\\mu$eV ] & [ $\\mu$eV ] & [meV] & $ [\\%] $ & $ [\\%] $ \\\\\n \\hline \n GaN & 400 & 3.954 & $\\pm$ 13.24 & $\\pm$ 52.79 & 11.97 & 0.22 & 0.88 \\\\\n \\hline\n ZnO & 320 & 3.855 & $\\pm$ 14.67 & $\\pm$ 58.39 & 14.65 & 0.20 & 0.80 \\\\\n \\hline\n InN & 345 & 4.885 & $\\pm$ 14.71 & $\\pm$ 59.49 & 16.20 & 0.18 & 0.73 \\\\\n \\hline \n AlN & 6.45 & 3.767 & $\\pm$ 3.98 & $\\pm$ 15.81 & 70.86 & 0.011 & 0.15 \\\\ \n \\hline\n\\end{tabular}\n\\end{center}\n\\caption{Intrinsic spin-orbit interaction corrections to the $E_2^*$ energy level of hydrogenic \ndonor impurities in various semiconductors with wurtzite crystal structure.\nThe parameters $\\gamma$ and $b$ are also indicated.\nThe effective masses and dielectric constants used to calculate $E_2^{*}$ are given in Table II.}\n\\end{table}\n\n\\section{Derivation of the Extrinsic spin-orbit interaction}\n\\label{cap:wurtzita-ext}\n\nThe spin-orbit Hamiltonian of an electron in vacuum in the presence of an electrostatic\npotential $V_0(\\mathbf{r})$ is given by\n\\begin{equation}\n H_{\\text{so}} = \\lambda \\,\\boldsymbol{\\hat{\\sigma}} \\cdot \n \\mathbf{k}\\times \\nabla V_{0}(\\mathbf{r}).\n\\end{equation}\nWhen the electron is immersed in a semiconductor in the presence of a mesoscopic potential \n$V(\\mathbf{r})$, the effective extrinsic spin-orbit Hamiltonian takes the form \n\\begin{equation}\n H_{\\text{ext}} = \\lambda^{\\ast} \\, \\boldsymbol{\\hat{\\sigma}} \\cdot \n \\mathbf{k} \\times \\boldsymbol{\\nabla}V(\\mathbf{r}),\n\\label{eq:HSOeff}\n\\end{equation}\nwhere $\\lambda^{\\ast}$ is an effective coupling constant.\nThis expression is valid for semiconductors with zincblende crystal structure,\nwhich presents a basic cubic symmetry.\nThe wurtzite crystal structure has less symmetry than the zincblende, due to the special role of\nits c-axis.\nThis lack of isotropy is also present in the intrinsic spin-orbit coupling given above, \nEq.\\ \\eqref{eq:int_completa}.\nIn what follows we shall derive an expression analogous to Eq.\\ \\eqref{eq:HSOeff} for semiconductors \nwith wurtzite crystal structure.\n\n\nWe start with the $\\mathbf{k} \\cdot \\mathbf{p}$ crystal Hamiltonian:\n\\begin{equation}\n\\label{eq:kdotp}\n \\mathcal{H} = \\mathcal{H}_0 + \\mathcal{H}_{\\mathbf{k} \\cdot \\mathbf{p}} + \\mathcal{H}_{\\text{so}}\n\\end{equation}\nwhere \n\\begin{eqnarray}\n \\mathcal{H}_0 &=& \\frac{P^2}{2 m} + U(r) , \\nonumber \\\\\n \\mathcal{H}_{\\mathbf{k} \\cdot \\mathbf{p}} &=& \\frac{\\hbar}{m}\\mathbf{k}\\cdot \\mathbf{p} + \n \\frac{\\hbar^{2} k^{2}}{2\\,m} , \\nonumber \\\\\n \\mathcal{H}_{\\text{so}} &=& \\frac{\\lambda}{\\hbar} \\, \\boldsymbol{\\hat{\\sigma}} \\cdot \n \\mathbf{p} \\times \\boldsymbol{\\nabla}U ,\n\\end{eqnarray}\nand $U$ is the periodic crystal potential.\nWe will write the matrix of $\\mathcal{H}$ in the common basis of $H_0$ and $J_z$ given by:\n\\begin{eqnarray}\n && v_1 = \\arrowvert i S \\uparrow \\rangle \\nonumber \\\\\n && v_2 = \\arrowvert i S \\downarrow \\rangle \\nonumber \\\\\n && v_3 = -\\frac{1}{\\sqrt2} \\left \\arrowvert (X+iY) \\uparrow \\right\\rangle \\nonumber \\\\\n && v_4 = -\\frac{1}{\\sqrt6} \\Big[ \\arrowvert (X+iY)\n \\downarrow\\rangle -2\\arrowvert Z \\uparrow \\rangle \\Big] \\nonumber \\\\\n && v_5 = \\frac{1}{\\sqrt6} \\Big[ \\arrowvert (X-iY) \\uparrow\\rangle \n +2\\arrowvert Z \\downarrow \\rangle \\Big] \\nonumber \\\\\n && v_6 = \\frac{1}{\\sqrt2}\\arrowvert (X-iY) \\downarrow\\rangle \\nonumber \\\\\n && v_7 = -\\frac{1}{\\sqrt3} \\Big[ \\arrowvert (X+iY) \\downarrow\\rangle\n +\\arrowvert Z \\uparrow \\rangle \\Big] \\nonumber \\\\\n && v_8 = -\\frac{1}{\\sqrt3}\\Big[ \\arrowvert (X-iY) \\uparrow\\rangle \n -\\arrowvert Z \\downarrow \\rangle \\Big] . \n\\end{eqnarray}\nHere $\\arrowvert S \\eta \\rangle$ are conduction-band $s$-states, with energy $E_c$, and \n$\\arrowvert X \\eta \\rangle$, $\\arrowvert Y \\eta \\rangle$, and $\\arrowvert Z \\eta \\rangle$ \nare valence-band $p$-type states, with energy $E_v$.\nThe energy gap is given by $E_g = E_c-E_v$.\nWe calculate the matrix elements \n$\\mathcal{H}_{ij} = \\langle v_i\\,\\arrowvert \\mathcal{H} \\arrowvert v_j\\,\\rangle$,\nwhere $ \\{i,j = 1, \\ldots, 8 \\} $, and obtain\n\n\n\\begin{widetext}\n\\begin{equation}\n \\mathcal{H} = \\left(\\begin{matrix} \n E_c & 0 & \\frac{-1}{\\sqrt2}P_2k_+ & \\sqrt{\\frac{2}{3}} P_1k_z & \\frac{1}{\\sqrt6}P_2k_- & 0 & \\frac{-1}{\\sqrt3}P_1k_z & \\frac{-1}{\\sqrt3}P_2k_- \\\\\n \n 0 & E_c & 0 & \\frac{-1}{\\sqrt6}P_2k_+ & \\sqrt{\\frac{2}{3}} P_1k_z & \\frac{1}{\\sqrt2}P_2k_- & \\frac{-1}{\\sqrt3}P_2k_+ & \\frac{1}{\\sqrt3}P_1k_z \\\\\n \n \\frac{-1}{\\sqrt2}P_2k_- & 0 & E_v & 0 & 0 & 0 & 0 & 0 \\\\\n \n \\sqrt{\\frac{2}{3}} P_1k_z & \\frac{-1}{\\sqrt6}P_2k_- & 0 & E_v & 0 & 0 & 0 & 0 \\\\\n \n \\frac{1}{\\sqrt6}P_2k_+ & \\sqrt{\\frac{2}{3}} P_1k_z & 0 & 0 & E_v & 0 & 0 \\\\\n \n 0 & \\frac{1}{\\sqrt2}P_2k_+ & 0 & 0 & 0 & E_v & 0 & 0 \\\\\n \n \\frac{-1}{\\sqrt3}P_1k_z & \\frac{-1}{\\sqrt3}P_2k_- & 0 & 0 & 0 & 0 & E_v- \\Delta_0 & 0 \\\\\n \n \\frac{-1}{\\sqrt3}P_2k_+ & \\frac{1}{\\sqrt3}P_1k_z & 0 & 0 & 0 & 0 & 0 & E_v -\\Delta_0 \\\\\n \\end{matrix}\\right)\n\\end{equation}\n\\end{widetext}\nwhere $k_{\\pm} = k_x \\pm ik_y$ and \n$\\Delta_0 = \\frac{\\hbar}{4 m^2 c^4} \\langle X | \\frac{\\partial U}{\\partial x} P_y - \n\\frac{\\partial U}{\\partial y} P_x | Y \\rangle$ is the spin-orbit splitting of the valence bands.\nWe have defined the constants $P_1$ and $P_2$ coming from the matrix elements:\n\\begin{equation}\n \\frac{\\hbar}{m}\\langle -iS\\,\\downarrow\\arrowvert\\mathbf{k} \\cdot \\mathbf{p} \n \\arrowvert Z\\,\\downarrow\\rangle\\,= \\,-i \\frac{\\hbar}{m} k_z \n \\langle S\\,\\arrowvert p_z \\arrowvert Z\\rangle \\equiv k_z P_1 , \\nonumber\n\\end{equation}\n\\begin{equation}\n \\frac{\\hbar}{m}\\langle -iS\\,\\downarrow\\arrowvert\\mathbf{k} \\cdot \\mathbf{p} \n \\arrowvert X\\,\\downarrow\\rangle\\,= \\, -i \\frac{\\hbar}{m} k_x \n \\langle S\\,\\arrowvert p_x \\arrowvert X\\rangle \\equiv k_x P_2 .\n\\end{equation}\n\n\nWe now introduce the impurity potential $V(\\mathbf{r})$, which varies slowly in the length \nscale of the lattice constant.\nIts matrix elements in the basis \\{$v_i$\\} are essentially diagonal thanks to the orthogonality \nof the basis set and its slow variation in atomic scale.\nIn short, we are applying here the envelope function approximation.\nThe matrix of $\\mathcal{H}+V$ can be expressed in a compact form using the matrices $\\mathbf{T}$ \nfamiliar from group theory:\\cite{win}\n\\begin{equation}\n T_x = \\frac{1}{3\\sqrt{2}} \n \\left(\\begin{matrix} \n -\\sqrt{3} & 0 & 1 & 0\\\\\n 0 & -1 & 0 & \\sqrt{3} \n \\end{matrix}\\right),\n\\end{equation}\n\\begin{equation}\n T_y = \\frac{-i}{3\\sqrt{2}} \n \\left(\\begin{matrix} \n \\sqrt{3} & 0 & 1 & 0\\\\\n 0 & 1 & 0 & \\sqrt{3} \n \\end{matrix}\\right),\n\\end{equation}\n\\begin{equation}\n T_z = \\frac{\\sqrt{2}}{3} \n \\left(\\begin{matrix} \n 0 & 1 & 0 & 0 \\\\\n 0 & 0 & 1 & 0 \n \\end{matrix}\\right).\n\\end{equation}\nUsing these matrices, the Hamiltonian matrix becomes:\n\\begin{equation}\n \\left(\\begin{matrix} \n (E_c+V) \\mathbb{I}_{2\\times 2} & \\sqrt3 P_1\\mathbf{T} \\cdot \\mathbf{k}_{\\alpha} & \n \\frac{-1}{\\sqrt3}P_1 \\boldsymbol{\\hat{\\sigma}} \\cdot \\mathbf{k_{\\alpha}} \\\\\n \n \\sqrt3 P_2\\mathbf{T}^{\\dagger} \\cdot \\mathbf{k}_{\\alpha} & (E_v +V) \\mathbb{I}_{4\\times 4} & 0 \\\\\n \n \\frac{-1}{\\sqrt3}P_1 \\boldsymbol{\\hat{\\sigma}} \\cdot \\mathbf{k_{\\alpha}} & 0 & \n (E_v -\\Delta_0+V) \\mathbb{I}_{2\\times 2} \n \\end{matrix}\\right),\n\\end{equation}\nwhere $\\mathbf{k}_{\\alpha}=(\\alpha k_x, \\alpha k_y, k_z)$ and $\\alpha = P_2\/P_1$.\n\n\nFollowing the application of the Foldy-Wouthuysen transformation described by Winkler for \nzincblende semiconductors,\\cite{win} we obtain an effective equation, restricted to the\nconduction band, for the electronic states in the donor impurity\n\\begin{widetext}\n\\begin{equation}\n \\left[ \\mathbf{T} \\cdot \\mathbf{k}_{\\alpha} \\frac{3P_1^{2}}{E-V+E_g}\\mathbf{T}^{\\dagger} \n \\cdot \\mathbf{k}_{\\alpha} + \\boldsymbol{\\hat{\\sigma}} \\cdot \\mathbf{k_{\\alpha}}\n \\frac{P^{2}_1}{3(E-V+E_g)}\\boldsymbol{\\hat{\\sigma}} \\cdot \\mathbf{k_{\\alpha}} \\right]\\psi_c=(E-V)\\psi_c .\n\\end{equation}\n\\end{widetext}\nUsing the relation \n$ (\\boldsymbol{\\hat{\\sigma}} \\cdot \\mathbf{A}) \n(\\boldsymbol{\\hat{\\sigma}} \\cdot \\mathbf{B}) = \\mathbf{A}\\cdot \\mathbf{B} + i\\boldsymbol{\\hat{\\sigma}}\n\\cdot \\left(\\mathbf{A} \\times \\mathbf{B} \\right)$, \nwe obtain two terms from the second term in the above equation, one of which corresponds to the effective\nspin-orbit interaction in the conduction band:\n\\begin{equation}\n H_{\\text{ext}} = \\lambda^{*}_w \\, \n \\boldsymbol{\\hat{\\sigma}} \\cdot \\left(\\mathbf{k_{\\alpha}} \\times \\boldsymbol{\\nabla}_{\\alpha} V \\right).\n\\label{eq:Hext}\n\\end{equation}\nwhere we defined $\\boldsymbol{\\nabla}_{\\alpha} \\equiv (\\alpha \\frac{\\partial}{\\partial x}, \n\\alpha \\frac{\\partial}{\\partial y}, \\frac{\\partial}{\\partial z} )$.\nWe have identified the coupling constant for the extrinsic spin-orbit interaction in wurtzite semiconductors:\n\\begin{equation}\n \\lambda^{*}_w = \\frac{\\epsilon P_{1}^2}{3}\\left[\\frac{2}{E^{2}_g} - \n \\frac{1}{(E_g +\\Delta_0)^{2}} \\right],\n\\label{eq:lambda}\n\\end{equation}\nanalogous to the known coupling constant $\\lambda^{*}$ in Eq.\\ \\eqref{eq:HSOeff} for zincblende materials.\n\n\nNote the factor $\\alpha \\equiv P_2\/P_1$ in Eq.\\ \\eqref{eq:Hext}, which reflects the anisotropy\nof the wurtzite crystal structure.\nThe Coulomb potential of the hydrogenic impurity, $V(\\mathbf{r})$, that appears in Eq.\\ \\eqref{eq:Hext},\nwas introduced in Eq.\\ \\eqref{eq:Hzero}.\nActually, the spherically symmetric form given after Eq.\\ \\eqref{eq:Hzero} is a simplified expression \nwhich does not include the effect of the anisotropic effective mass and dielectric constant of wurtzite\ncrystal structures.\\cite{han}\nAs a first approximation, here we will work with this spherically symmetric Coulomb potential\nand will also disregard the $\\alpha$-dependence of $\\boldsymbol{\\nabla}_{\\alpha}$ and $\\mathbf{k_{\\alpha}}$.\nA complete treatment of the anisotropy effects would require considering the modified\neigenvalue problem of the anisotropic hydrogenic impurity, and then the effect of the\nfactor $\\alpha$ in the spin-orbit interaction.\nWe leave this refined treatment for future work.\nIn Table II we present the values of $\\lambda^*_w$ for GaN, ZnO, InN, and AlN,\nalong with the material parameters needed to evaluate Eq.\\ \\eqref{eq:lambda}.\n\n\n\\begin{table}[ht]\n\\label{tab:lambda}\n\\begin{center}\n\\begin{tabular} {| l | c | c | c | c | c |}\n\\hline\n & $ m^*\/m_0$ & $\\epsilon$ & $ E_g $ & $\\Delta_0 $ & $\\lambda^{*}_w $\\\\\n & & & [eV] & [meV] & [$10^{-2}$ \\AA$^2$] \\\\\n \\hline \n GaN & 0.32 & 9.53 & 3.51 & 72.9 & 5.95 \\\\\n \\hline\n ZnO & 0.32 & 8.62 & 3.44 & 43 & 3.08 \\\\\n \\hline\n InN & 0.26 & 7.39 & 0.78 & 40 & 1.33 \\\\\n \\hline \n AlN & 0.38 & 4.27 & 5.4 & -58.5 & -1.04 \\\\ \n \\hline\n\\end{tabular}\n\\end{center}\n\\caption{Coupling constant of the effective extrinsic spin-orbit interaction \n and auxiliary material parameters\\cite{han} for wurtzite semiconductors.}\n\\end{table}\n\n\n\n\\section{Extrinsic spin-orbit corrections to the $2p$ level of hydrogenic impurities}\n\nUsing Eqs.\\ \\eqref{eq:Hext} and \\eqref{eq:lambda} we obtain the Hamiltonian of the extrinsic spin-orbit \ninteraction due to the Coulomb potential of the hydrogenic the donor impurity:\n\\begin{equation}\n H_{\\text{ext}} = \\frac{24a^{*3}}{r^3} \\xi_{2p}^{\\ast} \\, \\mathbf{L} \\cdot \\mathbf{S} ,\n\\label{eq:ext_final}\n\\end{equation}\nwhere\n\\begin{equation}\n \\xi^{\\ast}_{2p} \\equiv \\frac{e^2 \\lambda^{\\ast}_w}{24 \\epsilon \\hbar^2 a^{*3}} .\n\\end{equation}\nAs anticipated above, we simplified the Hamiltonian $H_{\\text{ext}}$ by setting the \nratio $\\alpha=1$ in Eq.\\ \\eqref{eq:Hext}.\nWe thus revert to the standard spin-orbit coupling of the hydrogen atom, but take into account\nthe appropriate, renormalized, coupling constant $\\lambda^*_w$.\nThe calculation of the first-order corrections to the $2p$ energy level of the impurity due\nto $H_{\\text{ext}}$ now follows the standard treatment of spin-orbit interaction in the\nhydrogen atom.\nThe common eigenvalues of $H_{\\text{ext}}$ and $J^2$ are given by:\n\\begin{equation}\n \\varepsilon_1 = \\frac{1}{2}\\xi_{2p}^{\\ast} \\Big[ \\frac{3}{4}-2-\\frac{3}{4} \\Big]\\hbar^2 = \n - \\xi_{2p}^{\\ast} \\hbar^2\n\\end{equation}\nfor $j=1\/2$, and\n\\begin{equation}\n \\varepsilon_2 = \\frac{1}{2} \\xi_{2p}^{\\ast} \\Big[ \\frac{15}{4}-2-\\frac{3}{4} \\Big]\\hbar^2 = \n \\frac{1}{2} \\xi_{2p}^{\\ast} \\hbar^2\n\\end{equation}\nfor $j=3\/2$. \nWe thus obtain for the energy corrections:\n\\begin{eqnarray}\n && \\varepsilon_1 = - \\frac{e^2 \\lambda^{\\ast}}{12 \\epsilon a^{\\ast 3}} \\equiv -2 \\,\\beta \\nonumber \\\\\n && \\varepsilon_2 = \\frac{e^2 \\lambda^{\\ast}}{24 \\epsilon a^{\\ast 3}} \\equiv \\beta . \n \\label{eq:bett}\n\\end{eqnarray}\nThe numerical values of $\\varepsilon_1$ and $\\varepsilon_2$ are shown in Table III,\ntogether with the parameters needed for their evaluation.\nWe also give the energy variation as a percentage of the unperturbed energy,\n$(\\varepsilon_2 - \\varepsilon_1) \/ E_2^{\\ast}$.\nOne can see that the splitting due to extrinsic spin-orbit interaction is four orders of magnitude\nsmaller than the energy of the original level.\nThis ratio is small but it is not negligible as it is, in fact, one order of magnitude larger than \nthe one obtained for the hydrogen atom, which is equal to 0.00133 \\%.\nIt should be emphasized that this comparison between the hydrogenic impurity and the hydrogen atom\nwas not obvious a priori, since the renormalization of the coupling constant $\\lambda$ is very \npronounced (6 orders of magnitude) and could have produced results radically different.\n\nWe remark that for aluminum nitride (AlN) the relation between $\\varepsilon_1$ and $\\varepsilon_2$ \nis inverted. \nThis peculiarity originates in the particular characteristics of its electronic structure, \nwhich cause $\\lambda^{\\ast}$ to become negative. \n\n\\begin{table}[ht]\n\\label{tab:energias_SO_ext}\n\\begin{center}\n\\begin{tabular} {| l | c | c | c | c | c| c |}\n\\hline\n & $ a^{\\ast} $ & $ \\lambda^{\\ast}_w$ & $ \\text{Ry}^{\\ast} $ & $ \\varepsilon_1 $ & $ \\varepsilon_2 $ & $ (\\varepsilon_2 - \\varepsilon_1)\/E_2^{\\ast} $ \\\\\n & [\\AA] &[$10^{-2}$ \\AA$^2$]& [meV] & [$\\mu$eV] & [$\\mu$eV] & $ [\\%] $ \\\\\n \\hline \n GaN & 15.8 & 5.95 & 11.97 & -1.94 & 0.968 & 0.024 \\\\\n \\hline\n ZnO & 14.1 & 3.08 & 14.65 & -1.56 & 0.779 & 0.016 \\\\\n \\hline\n InN & 15.2 & 1.33 & 16.20 & -0.625 & 0.313 & 0.0058 \\\\\n \\hline \n AlN & 5.9 & -1.04 & 70.85 & 14.50 & -7.25 & -0.031 \\\\ \n \\hline\n\\end{tabular}\n\\end{center}\n\\caption{\\footnotesize{Extrinsic spin-orbit corrections to the $2p$ energy level of hydrogenic \ndonor impurities for four important wurtzite semiconductores, along with relevant material \nparameters.}}\n\\end{table}\n\n\n\\section{Combined intrinsic and extrinsic spin-orbit interactions}\n\nWe now study the effects of the intrinsic and extrinsic spin-orbit interactions combined on \nthe $E_2^*$ energy level of hydrogenic impurities in wurtzite semiconductors.\nThus, we now consider the complete Hamiltonian, Eq.\\ \\eqref{eq:Htot}, with $H_{\\text{int}}$\nand $H_{\\text{ext}}$ given in Eqs.\\ \\eqref{eq:int_completa} and Eq.\\ \\eqref{eq:ext_final},\nrespectively.\nWe will perform again a first-order perturbative treatment, now considering the full spin-orbit\nHamiltonian $H_{\\text{int}}+H_{\\text{ext}}$ as the perturbation.\nTo that end, we will express the spin-orbit coupling in the so-called uncoupled basis states \nof the $E_2^*$ shell, that is, $\\left\\lbrace|l, s; m_l,m_s\\rangle\\right\\rbrace$, used in \nSection \\ref{sec:efa-impurity} to treat the intrinsic-alone case.\nThe matrix of the combined spin-orbit Hamiltonian in this basis is:\n\\begin{widetext}\n\\begin{equation}\n H_{\\text{ext}} + H_{\\text{int}} = \\left(\\begin{matrix} \n -\\varepsilon & 0 & 0 & A & 0 & 0 & 0 & 0 \\\\\n \n 0 & -\\varepsilon & 0 & 0 & 0 & 0 & B & 0 \\\\\n \n 0 & 0 & \\beta-\\varepsilon& 0 & 0 & 0 & 0 & 0 \\\\\n \n A & 0 & 0 & -\\beta-\\varepsilon & \\sqrt2\\,\\beta & 0 & 0 & 0 \\\\\n \n 0 & 0 & 0 & \\sqrt2\\,\\beta& -\\varepsilon & 0 & 0 & 0 \\\\\n \n 0 & 0 & 0 & 0 & 0 & -\\varepsilon & \\sqrt2\\,\\beta & 0 \\\\\n \n 0 & B & 0 & 0 & 0 & \\sqrt2\\,\\beta & -\\beta-\\varepsilon & 0 \\\\\n \n 0 & 0 & 0 & 0 & 0 & 0 & 0 & \\beta-\\varepsilon \\\\\n \\end{matrix}\\right) ,\n\\end{equation}\n\\end{widetext}\nwhere $\\beta$ was defined in Eq.\\ \\eqref{eq:bett} and $A$ and $B$ have been defined in \nEqs.\\ \\eqref{eq:AB}. \n\nThe characteristic polynomial that solves the eigenvalue problem is:\n\\begin{equation}\n \\varepsilon^2 (\\varepsilon-\\beta)^2(\\varepsilon\\,\\beta+\\varepsilon^2-A^2-2\\,\\beta^2)\n (\\varepsilon\\,\\beta+\\varepsilon^2-B^2-2\\,\\beta^2)=0 ,\n\\end{equation}\nand the corresponding eigenvalues are:\n\\begin{equation}\n \\varepsilon^2 = 0 \\,\\,\\,\\,\\,\\,\\,\\,\\,\\,\\,\\, \\Rightarrow \\,\\,\\,\\,\\, \\varepsilon_{1,2} = 0, \n \\label{eq:rot}\n\\end{equation}\n\\begin{equation}\n \\,\\,\\,(\\varepsilon-\\beta)^2 = 0 \\,\\, \\Rightarrow \\,\\,\\,\\,\\,\\,\\, \\varepsilon_{3,4} = \\beta ,\n \\label{eq:rot1}\n\\end{equation}\n\\begin{equation}\n \\varepsilon^2+\\varepsilon\\beta-(A^2+2\\beta^2) = 0 \\,\\,\\,\\,\\,\\,\\,\\, \\Rightarrow \\,\\,\\,\\,\\, \\varepsilon_{5,6} = \n -\\frac{\\beta}{2} \\pm \\sqrt{\\frac{9}{4}\\beta^2 + A^2}, \n \\label{eq:rot2}\n\\end{equation}\n\\begin{equation}\n \\varepsilon^2+\\varepsilon\\beta-(B^2+2\\beta^2) = 0 \\,\\,\\,\\,\\,\\, \\Rightarrow \\,\\,\\,\\,\\, \\varepsilon_{7,8} = \n -\\frac{\\beta}{2} \\pm \\sqrt{ \\frac{9}{4} \\beta^2 + B^2}.\n \\label{eq:rot3}\n\\end{equation}\nWe remark that the energy corrections given in Eqs.\\ \\eqref{eq:rot}--\\eqref{eq:rot3} contain the \nprevious cases (intrinsic and extrinsic spin-orbit interactions acting alone) in the appropriate \nlimits, and they are represented schematically in Fig. \\ref{fig:combinado}.\nThis plot offers a qualitative view of the energy splittings and shows the greater breaking of\ndegeneracy caused by the combined action of both spin-orbit couplings.\nFinally, we calculate the eigenvalues $\\varepsilon_i$ ($i=3, \\ldots ,8$) using the values \nof $A$, $B$ y $\\beta$ corresponding to GaN, ZnO, InN y AlN; the results are given in Table IV.\n\n\\begin{figure}\n \\centerline{\\includegraphics[scale=0.6]{figure1.pdf}}\n \\caption{\\footnotesize{Schematic representation of the corrections to the energy level $E_2^*$ \n of hydrogenic donor impurities to first order in the intrinsic and extrinsic spin-orbit \n interactions combined.}\n \\label{fig:combinado}}\n\\end{figure}\n\n\n\\begin{table}[ht]\n\\label{tab:comb}\n\\begin{center}\n\\begin{tabular} {| l | c | c | c | c | c |}\n\\hline\n & $\\varepsilon_{3,4}$ & $\\varepsilon_5$ & $\\varepsilon_6$ & $\\varepsilon_7$ & $\\varepsilon_8$ \\\\\n & [$\\mu$eV] & [$\\mu$eV] & [$\\mu$eV] & [$\\mu$eV] & [$\\mu$eV] \\\\\n \\hline \n GaN & 0.23 & -13.36 & 13.13 & -54.71 & 54.48 \\\\\n \\hline\n ZnO & 4.90 & -49.39 & 44.94 & -60.78 & 56.33 \\\\\n \\hline\n InN & 0.76 & -15.42 & 14.66 & -61.31 & 60.63 \\\\\n \\hline \n AlN & -20.15 & -20.44 & 40.58 & -17.56 & 37.70 \\\\ \n \\hline\n\\end{tabular}\n\n\\caption{\\footnotesize{Energy corrections due to the combined intrinsic and extrinsic spin-orbit\ninteractions to the energy level $E_2^*$ of hydrogenic donor impurities for four important \nbinary semiconductors with wurtzite crystal structure.}}\n\\end{center}\n\\end{table}\n\n\n\\section{Conclusion}\n\\label{sec:conclusion}\n\nWe have studied theoretically the effects of the spin-orbit interaction on the\n$E_2^*$ energy level of hydrogenic donor impurities embedded in semiconductors\nwith wurtzite crystal structure.\nBoth the intrinsic (Dresselhaus) and extrinsic (Rashba) spin-orbit interations\nhave been considered, first acting separately and then together.\nThe study was carried out at the level of first-order perturbation theory,\nwhich turns out to be appropriate given the relative magnitude of the corrections\nto the unperturbed energy spacings.\nFurthermore, in order to evaluate the extrinsic spin-orbit interaction \nit was necessary to calculate the renormalized coupling constant $\\lambda^*$ \nfor wurtzite semiconductors from 8-band Kane theory.\n\n\nWe applied our calculations to four currently important semiconductors, i.e.\\ GaN, ZnO \nInN, and AlN.\nA general conclusion of these calculations is that both spin-orbit couplings produce\nrelative energy corrections that are bigger than the standard spin-orbit corrections\nto the $E_2$ energy level of the hydrogen atom.\nWhile for GaN, ZnO, InN we conclude that the intrinsic spin-orbit interaction\nproduces larger energy corrections than the extrinsic one, that is not the case for AlN,\nwhere both interactions have comparable effects.\nAnother anomaly shown by AlN is the fact that its effective coupling constant \n$\\lambda^*$ is negative.\nThis causes the eigenvalues $\\varepsilon_{3,4}$, which are positive for GaN, ZnO \nand InN, to become negative for AlN.\nThese anomalies of AlN are due to the specific features of its electronic structure\nwhich determine the relevant parameters $\\Delta_0$, $\\gamma_w$ and $a^{\\ast}$.\nFinally, we have found that the combined action of both types of spin-orbit coupling \nleads to an almost complete breaking of the degeneracy of the unperturbed energy level.\n\n\\bigskip\n\n \n\\acknowledgments\nWe thank Rodolfo Jalabert and Dietmar Weinmann for discussions that motivated this work.\nWe gratefully acknowledge financial support from Projects UBACYT and CONICET PIP.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\\subsection{Motivation}\n\\input{sections\/intro\/motivation.tex}\n\\subsection{Related Work and Main Contributions}\n\\input{sections\/intro\/contributions.tex}\n\\subsection{Notation}\n\\input{sections\/intro\/notation.tex}\n\n\\section{Preliminaries} \\label{sec:prelimiaries}\n\\subsection{System Model}\\label{sec:sys_model}\n\\input{sections\/preliminaries\/sys_model.tex}\n\\subsection{Metrics}\n\\input{sections\/preliminaries\/metrics.tex}\n\\section{Jammer Model}\n\\input{sections\/preliminaries\/jammer_model.tex}\n\n\\section{Resilient Receiver Design} \\label{sec:transceiver}\n\\subsection{General Considerations}\\label{sec:look_sinr}\n\\input{sections\/resilient_design\/look_sinr.tex}\n\\input{sections\/resilient_design\/separate_design.tex}\n\n\n\\section{Numerical Results} \\label{sec:results}\n\\subsection{Methodology}\n\\input{sections\/results\/methodology.tex}\n\\subsection{Simulation Studies}\n\\input{sections\/results\/simulation.tex}\n\n\\section{Conclusions} \\label{sec:conclusions}\n\\input{sections\/conclusions.tex}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{Closed-Form Design}\\label{subsec:closed_form}\nThe receivers' goal is to ensure reliable communication for all uplink participants,\ni.e. we are interested in maximizing the individual per-stream SINRs $\\gamma_{k}^{B}$ or \nat least ensure a certain quality of service (QoS), e.g. the $\\gamma_{k}^{B}$'s lie above a given threshold $\\gamma_{0}$.\n\nThe maximization of the lower bound in Eq. \\eqref{eq:lb} is a Rayleigh quotient maximization problem with the\nstandard solution being given by:\n\\begin{align}\n \\label{eq:maxsinr}\n \\vb_{k} = \\br{\\bm{B}_{k} + \\eta\\bm{A}_B(\\bm{\\theta}_G)\\abg^{\\mathsf{H}}}^{-1}\\ch_{k}\\wa_{k}\n\\end{align} \nfor all $k = 1, \\dots, K$. \nThis is the most simple way to leverage AoA information at the receiver, since the solution is given in closed form\nand computationally tractable, since $\\bm{A}_B(\\bm{\\theta}_G)\\abg^{\\mathsf{H}}$ can be easily computed without matrix multiplications \nby only considering the phase shifts between the impinging jamming signals.\n\nNote, that the resulting filter still depends on $\\eta$, which in general is not available at the receiver.\nThe most easy way to address this problem is to consider $\\eta$ as a hyperparameter which controls how much the\njamming signals are suppressed. \nIndeed, this is justified by the fact that in the jammer-dominant regime, \nthe contribution of the noise variance $\\sigma^{2}$ is small compared to the one of the jammer, i.e\n$\\vb_{k}^{\\mathsf{H}} \\br{\\eta\\bm{A}_B(\\bm{\\theta}_G)\\abg^{\\mathsf{H}} + \\sigma^{2}\\bm{\\mathrm{I}}} \\vb_{k} \\approx \\eta\\vb_{k}^{\\mathsf{H}} \\bm{A}_B(\\bm{\\theta}_G)\\abg^{\\mathsf{H}} \\vb_{k}$.\nIn the following, we will propose other methods to overcome this dependency.\n\n\\subsection{Zero Forcing Design}\nAs it can be seen in Section \\ref{subsec:closed_form}, the main contribution to the jamming signals is given\nby the term \n\\begin{align}\n \\beta_{k} = \\normsq{\\bm{A}_B(\\bm{\\theta}_G)^{\\mathsf{H}} \\vb_{k}} \n\\end{align}\nwhich can be naturally interpreted as the receive beampattern of the receiver antenna array \\cite{van2004optimum}. \nThis indicates how much the signals coming from the jammer direction are emphasized, thus \nthe anti-jamming criterion can be formalized as a beampattern minimization problem in the same spirit as\n\\cite{mu_mimo_radar_com,su_pls_bf,jcas_sp_survey}.\nFurthermore, a natural way to ensure good per-stream SINRs for all users is to remove the interference,\nwhich can be realized with a zero-forcing (ZF) design. \nThus, the beampattern minimization problem under ZF constraint can be formulated analogously to \\cite{mu_mimo_radar_com} as\n\\begin{align}\n \\label{eq:bpmin_zf}\n \\min_{\\substack{\n \\vb_{k}\\\\\n k = 1, \\dots, K}\n }\\sum_{k=1}^{K}\\beta_{k} \\quad \\text{s.t.} \\quad \\bm{P}^{\\mathsf{H}} \\vb_{k} = \\bm{e}_{k} \\, \\forall k\n\\end{align}\nwhere we defined \n\\begin{equation}\n \\bm{P} = \\begin{bmatrix}\n \\bm{H}_{1}\\bm{w}_{1} & \\dots & \\bm{H}_{K}\\bm{w}_{K}\n \\end{bmatrix} \\in \\mathbb{C}^{N_{B} \\times K}\n\\end{equation}\nand $\\bm{e}_{k}$ is $k$-th standard basis vector in $\\mathbb{R}^{K}$.\nNote that this is a convex optimization problem easily solvable in polynomial time. \n\nFurthermore, if the matrix $\\bm{A}_B(\\bm{\\theta}_G)\\abg^{\\mathsf{H}}$ has full-rank, we can also give an analytic solution.\nLetting $\\bm{X} = \\bm{A}_B(\\bm{\\theta}_G)\\abg^{\\mathsf{H}}$, we construct the Lagrangian\n\\begin{align}\n \\mathcal{L}\\left(\\vb_{k}, \\bm{\\mu}_{k}\\right) = \n \\sum_{k=1}^{K} \\vb_{k}^{\\mathsf{H}}\\bm{X}\\vb_{k} + 2\\Re\\brcur{\\bm{\\mu}_{k}^{\\mathsf{H}} \\br{\\bm{P}^{\\mathsf{H}} \\vb_{k} - \\bm{e}_{k}}}\n\\end{align}\nwhere $\\bm{\\mu}_{k}\\in\\mathbb{C}^{K}$ are the Lagrange multipliers.\nSince the ZF-constraint is an equality constraint, the Karush-Kuhn-Tucker conditions are both necessary and sufficient.\nComputing the gradients wrt. $\\vb_{k}$ and $\\bm{\\mu}_{k}$, setting them to $\\bm{0}$, and solving for $\\vb_{k}$,\nwe obtain\n\\begin{align}\n \\label{eq:bpmin_zf_analytic}\n \\vb_{k} = \\bm{X}^{-1} \\bm{P} \\br{\\bm{P}^{\\mathsf{H}} \\bm{X}^{-1} \\bm{P}}^{-1} \\bm{e}_{k}\n\\end{align}\nBy standard linear algebra $\\bm{X}$ has full rank iff $L_{G}\\geq N_{B}$, which might pose a problem if the propagation channel\nof the jammer shows $L_{G} < N_{B}$ dominant paths. \nIn that case, one can easily ``pad'' the resolved AoAs $\\bm{\\theta_{G}}$ with \na set of different angles $\\{\\phi_{G, l}\\}_{l=1}^{N_{B}-L_{G}}$ thus enlarging the grid over which the beampattern should be minimized.\nThis idea can be naturally extended to the designs considered in Sections \\ref{subsec:closed_form} and \\ref{subsec:min_sinr}.\n\\subsection{Minimum SINR Design}\\label{subsec:min_sinr}\nAnother way to approach this problem is to minimize the influence of the jamming signals, \nwhile ensuring a minimum SINR $\\gamma_{0}$ in the jammer-free case.\nUsing this method one can address the disadvantages ZF-filters generally have in the low SNR regime, i.e. $\\sigma^{2} \\gg 0$.\nRewriting the beampattern $\\beta_{k}$ as\n\\begin{align}\n \\beta_{k} = \\normsq{\\bm{A}_B(\\bm{\\theta}_G)^{\\mathsf{H}} \\vb_{k}} = \\tr{\\bm{A}_B(\\bm{\\theta}_G)^{\\mathsf{H}} \\bm{V}_{k} \\bm{A}_B(\\bm{\\theta}_G)}\n\\end{align}\nwe cast the resulting problem as\n\\begin{align}\n \\label{eq:bpmin}\n \\min_{\n \\substack{\\bm{V}_{k} \\succcurlyeq \\bm{0}\\, \\bm{V}_{k} = \\bm{V}_{k}^{\\mathsf{H}} \\\\ \\text{rank}(\\bm{V}_{k}) = 1 \\\\\n k = 1, \\dots, K}\n } \\,\n \\sum_{k=1}^{K}\\beta_{k} \\quad \\text{s.t.} \\quad \\dfrac{\\tr{\\bm{A}_{k}\\bm{V}_{k}}}{\\tr{\\bm{B}_{k}\\bm{V}_{k}}} \\geq \\gamma_{0}\n\\end{align}\nwith $\\bm{V}_{k} = \\vb_{k} \\vb_{k} ^{\\mathsf{H}}$ as before and $\\bm{A_{k}}$ and $\\bm{B_{k}}$ as in equations \\eqref{eq:sh1} and \\eqref{eq:sh2}.\nNote that this problem is non-convex due to the rank-$1$ constraint on $\\bm{V}_{k}$. \nIn order to efficiently solve this problem, we apply semidefinite relaxation \\cite{sdr}, i.e. \nwe relax the rank-$1$ constraint, \nsolve the resulting convex problem and recover the solution via eigenvalue decomposition (EVD).\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzedut b/data_all_eng_slimpj/shuffled/split2/finalzzedut new file mode 100644 index 0000000000000000000000000000000000000000..ef9d2e166902fce07b52df15af31a71ae9d1a301 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzedut @@ -0,0 +1,5 @@ +{"text":"\\section{INTRODUCTION}\n\nThe first step of planet formation in protoplanetary disks is the collisional\ngrowth of (sub-)micron-sized dust particles (or aggregates), and consequent\nsedimentation. Dust settling and the formation of a dust layer at the midplane of\nthe disk play an important role in the subsequent planetesimal formation\nprocess. Enhancement of the particle density in the dust layer accelerates the\ncollisional growth rate. If the density enhancement is high enough,\nplanetesimals may form through the gravitational instability of the dust layer\n(Goldreich \\& Ward 1973; Sekiya 1983). However, the turbulent motion of the gas\nhinders the dust layer from thinning. Dust particles are stirred by\nturbulence and diffuse to high altitudes from the midplane. Thus, turbulence \nis an impediment to planetesimal formation. Turbulence excited by\nmagneto-rotational instability (MRI) is so strong that the dust layer cannot\nbecome thin enough to induce planetesimal formation through gravitational\ninstability (Johansen \\& Klahr 2005; Fromang \\&\\ Papaloizou 2006; Turner et\nal. 2006; Carballido et al. 2006, 2011; Fromang \\& Nelson 2009). In addition, the\nturbulence-induced collisional velocity can be high enough to destroy dust\naggregates (Carballido et al. 2008, 2010), though several mechanisms have been\nproposed to overcome this difficulty (Johansen et al. 2007, 2011; Lyra et al.\n2008, 2009). Even if the gas disk exhibits an initially laminar flow (such as is expected\nin the dead zone where the ionization fraction of the gas is too low to couple\nto the magnetic field, as mentioned in Gammie 1996 and Sano et al. 2000), dust \nsettling itself induces turbulence. The dust particles tend to rotate around \nthe central star faster than the gas because the gas experiences a pressure gradient \nforce acting in the opposite direction of the gravity of the star. When the dust layer \nthins and the dust-to-gas ratio in the layer approaches unity, the velocity\ndifference from that of the upper gas layer induces Kelvin-Helmholtz\n(KH) instability and excites turbulence in the dust layer (e.g.,\nGoldreich \\& Ward 1973; Sekiya \\& Ishitsu 2000, 2001; Garaud \\&\\ Lin\n2004). The velocity difference between the dust particles and the gas\ninside the dust layer also induces streaming instability (Youdin \\&\nGoodman 2005; Youdin \\& Johansen 2007; Johansen \\& Youdin 2007; see\nChiang \\& Youdin 2010 for review on the various instability in the dust\nlayer). Laminar gas disks are considered excellent sites for planetesimal\nformation. Thus, it is important to clarify how turbulence induced in\nthe dust layer diffuses the dust particles, versus dust\nsedimentation. Several analytical and numerical studies have focused on\nthis problem (Cuzzi et al. 1993; Champney et al. 1995; Sekiya 1998;\nDobrovolskis et al. 1999; Johansen et al. 2006; Michikoshi \\&\\ Inutsuka\n2006; Weidenschilling 2006, 2010; Bai \\&\\ Stone 2010a, 2010b).\n\nSekiya (1998, hereafter S98) analytically solved for the structure of the dust\nlayer under turbulence induced by KH instability. To derive the density\nprofile of the dust layer analytically, the author adopted several assumptions. First, the dust particles were assumed to be small enough and were coupled so tightly to the gas that\nthe dust and the gas could be treated as a single fluid. Second, the density\nstructure was adjusted to keep the dust layer marginally unstable\nto KH instability. Third, the effects of the Coriolis force and of the Keplerian\nshear were neglected. These assumptions allowed the author to determine\nanalytically the density structure of the dust layer and to discuss how far\nfrom the midplane the dust particles diffuse due to turbulence. The\nanalysis in S98 provided a useful guide for subsequent numerical studies\nthat investigated turbulence of the dust layer with more realistic assumptions.\nIn reality, however, dust particles that are not tightly coupled to gas\nplay an important role in inducing turbulence. Also, other kinds of instability,\nsuch as streaming instability, may occur before KH instability sets in (Bai \\&\nStone 2010a). Thus, in addition to performing numerical simulations, it is\ndesirable to include an analytic discussion on the structure of the dust layer using\na more general set of assumptions than adopted in S98.\n\nIn this paper, we revisit the analysis presented in S98 from a different point of view.\nIn S98, the density structure of the dust layer was determined by the condition\nthat the dust layer was marginally unstable to KH instability; i.e., the\nRichardson number, $\\mathrm{Ri}$, had a constant critical value. In the\nstability condition on $\\mathrm{Ri}$, the free energy due to the shear\nvelocity is compared with the energy needed to lift the material against the\nvertical gravity (Chandrasekhar 1961). In the case of a dust layer, the free\nenergy of the azimuthal velocity is compared with the gravitational potential in\nthe vertical direction. We also discuss the energetics of the dust layer to\ndetermine the density structure, but in this paper we focus on the\ngravitational potential in the radial direction. Sustaining steady turbulence\nrequires an energy supply. The source of the free energy for KH instability or\nstreaming instability is the velocity difference between the dust and the gas.\nBecause this free energy is consumed in the process of exciting turbulence, a velocity difference must be continually induced to provide steady turbulence. This velocity difference originates from the force balance in the radial direction\nbetween the stellar gravity, the centrifugal force, and the pressure gradient.\nIn other words, the velocity difference is induced because the gas resides in a slightly shallower effective potential (including both the\nstellar gravity and the gas pressure), than the effective potential for the dust. The dust\nloses its angular momentum due to gas drag and drifts towards the star, while the gas gains \nangular momentum and drifts outward. The total angular momentum remains\nconstant. Because the effective potential for the dust is deeper, this process releases \ngravitational energy that can be a source of the free energy for turbulence. (The relationship between the angular momentum variation, $\\Delta L$, and the energy\nvariation, $\\Delta E$, for circular orbiting material is $\\Delta E=\\Omega\\Delta\nL$, where the angular velocity, $\\Omega$, is different for the gas and for the\ndust.) If the gas disk is initially in a state of laminar flow, i.e., if gas\ndrag is the only process that exchanges angular momentum, then dust accretion is the source of the free energy used to induce turbulence. The accretion\nrate of the dust due to gas drag has been calculated in the literature\n(Nakagawa et al. 1986, hereafter NSH86; Weidenschilling 2003; Youdin \\& Chiang\n2004). Using this dust accretion rate, and assuming that a certain\nfraction of the accretion energy is transferred to turbulence, it is\nstraightforward to calculate the strength of the turbulence and to determine the\ndensity structure of the dust layer. In this paper, we show that an analysis of the\nenergetics based on the gravitational potential in the radial direction results in a dust layer structure that is qualitatively equivalent to the dust layer structure developed in\nthe analysis presented in S98. Most of the important properties of the S98 model\nare reproduced. Our analysis avoids some of the assumptions that were adopted in S98. \nSpecifically, we do not assume tight coupling between the gas and the dust, nor a\nmarginally unstable structure for KH instability. (The effects of the Coriolis\nforce and of Keplerian shear, which are still neglected in this paper, are\ndiscussed by Ishitsu \\& Sekiya 2002, 2003; G\\'omez \\& Ostriker 2005; Chiang\n2008; Barranco 2009; Lee et al. 2010.) Thus, the results of this paper may be applied to more general situations, including situations in which the\ndust is weakly coupled to the gas, or in which streaming instability acts as a source\nof turbulence.\n\nIn \\S \\ref{sec:assump}, we describe our model assumptions. In\n\\S \\ref{sec:eneinp}, the energy release rate of the accreting dust is\ncalculated. In \\S \\ref{sec:turbstr}, the strength of the turbulence is derived and\nits asymptotic forms in the limits of a small dust-to-gas ratio and a\ntight dust-gas coupling are discussed. In \\S \\ref{sec:comppre}, the results of\nour model are compared with the S98 model and with previous numerical\nsimulations of KH instability and streaming instability. In\n\\S \\ref{sec:dustvelocity}, we discuss how dust layer formation reduces\nthe radial drift velocity of the dust and the relative velocity of the dust\nparticles. In \\S \\ref{sec:ene-into-tur}, we check the applicability of our model to\nturbulence in the dust layer for the KH instability and streaming instability cases.\n\n\\section{DISK MODEL AND ASSUMPTIONS}\n\\label{sec:assump}\n\n\\subsection{Brief Model Description}\n\nWe consider a dust layer that forms after dust particles have settled to the\nmidplane of a gas disk around a star. Before the dust settles, the gas disk is\nassumed to be subject to laminar flow. This means that we consider the dead zone where MRI is\ninactive because of the low degree of ionization. Turbulent gas motion may be present even in the dead zone because sound waves propagate from the\nactive layers at high altitudes and disturb the gas in the dead zone (Fleming \\&\nStone 2003; Suzuki et al. 2010; Okuzumi \\& Hirose 2011). We assume that the\ndead zone is wide enough that the active layer cannot induce strong turbulence\nat the midplane.\n\nEven in an initially laminar disk, dust sedimentation and the formation of the\nmidplane dust layer cannot proceed in a perfectly undisturbed fashion. \nWhen the midplane dust-to-gas\nratio reaches a critical value, the velocity differences between the dust\nlayer and the upper gas layer, or between the individual dust particles and the\nsurrounding gas, will start to induce hydrodynamical instabilities, such as KH\ninstability or streaming instability (see Chiang \\& Youdin 2010 for review).\nConsequently, turbulent diffusion of the dust particles terminates further\ndust settling. It is expected that a steady dust layer forms, in which \nparticle settling and turbulent diffusion balance each other. In such a\nstate, vertical settling is no longer an energy source for turbulence,\nbut the radial accretion of the dust still proceeds. The turbulence is\nmaintained by the energy liberated from dust accretion to the star.\n\nWe calculate the dust accretion rate and consequent energy release, \nassuming that a steady dust layer has formed. The strength of the\nturbulence is then estimated. The balance between the turbulent diffusion\nand the dust settling determines the dust layer thickness. Thus, the structure of\nthe dust layer, the dust accretion rate, and the turbulence strength must be\ndetermined self-consistently.\n\nFor simplicity, we focus on a narrow axisymmetric ring region at a certain\nradius $r$ from the central star. We consider the dust layer structure only in\nthe vertical direction, neglecting any radial variation in the properties of\nthe gas disk and the dust layer. All the dust particles are assumed to have\na uniform size, i.e., the size distribution of the dust particles is neglected\nin this paper.\n\n\\subsection{Vertical Structure of the Dust Layer}\n\nIf the effect of gas drag on the dust were negligible, the dust particles\n would orbit at the Keplerian velocity \n$\\Omega_{\\mathrm{K}}=\\sqrt{(GM)\/r^{3}}$, where $G$ is the \ngravitational constant and $M$ is the mass of the\ncentral star. However, the orbital velocity of the gas is slightly less\nthan the Keplerian velocity because the direction of the gas pressure\ngradient is usually outward from the star. This cancels part of the \ncontribution to the orbital velocity from the gravity of the star. Given a negligible influence of \nthe drag force from the dust, the angular velocity of the gas would be\n\\begin{equation}\n\\Omega_{g}=\\sqrt{\\frac{GM}{r^{3}}(1-2\\eta)}\\approx\\Omega_{\\mathrm{K}}(1-\\eta).\n\\label{eq:omega_g\n\\end{equation}\nHere, $\\eta$ represents the deviation of the gas orbital velocity from the\nKeplerian value. It is half of the ratio between the gas pressure gradient\nforce and the gravity of the central star,\n\\begin{equation}\n\\eta=-\\frac{1}{{2\\rho_{g}r\\Omega_{\\mathrm{K}}^{2}}}\\frac{{\\partial P}}{{\\partial r}}\\ ,\n\\end{equation}\nwhere $\\rho_{g}$ and $P$ are the gas density and pressure, respectively.\n\nWe consider a thin dust layer in which most of the dust particles have settled\non the midplane of the gas disk. The gas disk is isothermal in the vertical\ndirection with a scale height of $h_{g}=c_{s}\/\\Omega_{\\mathrm{K}}$, where\n$c_{s}$ is the gas sound speed. Because the thickness of the dust layer is\nmuch smaller than the scale height of the gas disk, the gas density $\\rho_{g}$\nin the dust layer is treated as a constant. Its value is given by the column\ndensity of the gas disk $\\Sigma_{g}$ as\n\\begin{equation}\n\\rho_{g}=\\frac{\\Sigma_{g}}{\\sqrt{2 \\pi} h_{g}}\\ . \n\\label{eq:rhog\n\\end{equation}\nThe dust density profile in the vertical direction, $\\rho_{d}(z)$, is\ndetermined by the balance between the dust settling and turbulent diffusion\n(Youdin \\& Lithwick 2007, hereafter YL07)\n\\footnote{Note that the diffusion coefficient\ngiven by YL07 is derived assuming that the dust particles can be treated as\npassive particles. Problems may arise if this formula is used in a case \nwhere the dust layer is so dense that its dust-to-gas ratio is larger than unity. Although a\ncomparison of our model with the simulation by Johansen et al. (2006) shows \ngood agreement even for a dust-to-gas ratio as large as $40$ (see Fig.\n\\ref{fig:fdmid-joh}), modeling the effect of the dust inertia will be a \nsubject for future investigation using a refinement of the present\nmodel.}.\nThe dust particles have a\nstopping time, $\\tau_{\\mathrm{stop}}=T_{s}\\Omega_{\\mathrm{K}}^{-1}$, in which\nthe velocity difference from the background gas flow becomes $1\/e$ times due\nto gas drag, where $T_{s}$ is the non-dimensional stopping time. In a\nturbulent gas disk with a turbulent diffusion coefficient $D_g$,\nthe density profile of the dust layer composed of single-sized particles is\nGaussian, $\\rho_{d}=\\rho_{d,0}\\exp(-\\tilde{z}^{2})$, where $\\tilde{z\n=z\/(\\sqrt{2}h_{d})$ is the non-dimensional vertical coordinate normalized by\nthe dust scale height $h_d \\approx [D_g \/(\\Omega_{\\rm K} T_s)]^{1\/2}$.\nNote that the turbulent diffusion coefficient in the ``z-direction'',\n$D_g$, may have a significantly different value from the usual turbulent\nviscosity coefficient $\\nu_{\\rm acc}$ in accretion disks, where\n$\\nu_{\\rm acc}$ comes from the ``$r \\theta$-component'' of the Reynolds\nstress (e.g., Lesur \\& Ogilvie 2010). Nevertheless, we follow the conventional\n``$\\alpha$-prescription'', and express the diffusion coefficient as\n$D_g=\\alpha c_g h_g$ for simplicity. This prescription and equation (24)\nof YL07 give\n\\begin{equation}\nh_{d}=\\sqrt{\\frac{\\alpha}{{T_{s}}}}\\sqrt{\\frac{{1+T_{s}}}{{1+2T_{s}}}}h_{g}\\ .\n\\label{eq:hd\n\\end{equation}\nThe dust density profile is written as\n\\begin{equation}\n\\rho_{d}=f_{\\mathrm{mid}}\\rho_{g}\\exp(-\\tilde{z}^{2})\\ ,\n\\end{equation}\nwhere the midplane dust-to-gas ratio $f_{\\mathrm{mid}}$ is increased by a\nfactor $h_{g}\/h_{d}$ from the total dust-to-gas ratio or the\n``metallicity'' of the disk $Z_{\\mathrm{disk}}=\\Sigma_d \/ \\Sigma_g$,\n\\begin{equation}\nf_{\\mathrm{mid}}=Z_{\\mathrm{disk}}\\frac{{h_{g}}}{{h_{d}}}=Z_{\\mathrm{disk}}\\sqrt{\\frac{{T_{s}}}{\\alpha}}\\sqrt{\\frac{{1+2T_{s}}}{{1+T_{s}}}}\\ .\n\\label{eq:fdmid\n\\end{equation}\nWe define $\\beta(\\tilde{z})$ as the ratio of the total (dust and gas)\ndensity to the gas density,\n\\begin{equation}\n\\beta=\\frac{\\rho_{g}+\\rho_{d}}{\\rho_{g}}={1+f_{\\mathrm{mid}}\\exp(-\\tilde\n{z}^{2})}\\ . \n\\label{eq:beta\n\\end{equation}\n\n\\subsection{Typical Values of the Model Parameters}\n\nIn our model, the turbulent parameter $\\alpha$ and the structure of the dust\nlayer (i.e., the thickness, $h_{d}$, and the midplane dust-to-gas ratio,\n$f_{\\mathrm{mid}}$) are determined for a given parameter set (the stopping\ntime of the dust particles $T_{s}$, the disk metallicity\n$Z_{\\mathrm{disk}}$, and the parameters for the gas disk). The\nnon-dimensional variables of the result ($\\alpha$, $h_{d}\/r$, and\n$f_{\\mathrm{mid}}$) depend on the gas disk parameters only through\n\\begin{equation}\n\\tilde{\\eta}=\\eta^{2}\\left( \\frac{r}{h_{g}}\\right) ^{2}~, \n\\label{eq:etawave\n\\end{equation}\nwhere $\\tilde{\\eta}\\sim\\eta\\sim(h_{g}\/r)^{2}\\sim10^{-3}-10^{-2}$. The\nnumerical calculations in the subsequent sections mainly use a value\n$\\tilde{\\eta}=(0.05\/\\gamma)^{2}=9 \\times 10^{-4}$, where $\\gamma=5\/3$,\nin order to compare our model with the simulation presented by \nJohansen et al. (2006, hereafter JHK06). \nIn \\S\\ref{sec:BS10}, $\\tilde{\\eta}=0.05^{2}$ is used for comparison\nwith the results by Bai \\& Stone (2010), and in \\S\\ref{sec:dustvelocity},\n$\\tilde{\\eta}=2.92 \\times 10^{-3}$ is used to calculate the radial\ndrift velocity of the dust in the gas disk model by Hayashi (1981). The\n$\\tilde{\\eta}$ dependence of the results will be discussed in\n\\S\\ref{sec:disketa}. \n\n\n\\section{TURBULENCE ENERGY SUPPLY}\n\\label{sec:eneinp}\n\nThe energy supply for turbulence comes from liberation of the gravitational\nenergy of the dust. As the dust falls towards the central star, the dust\nparticles penetrate more deeply into the potential well of the star. Although the gas drifts\noutward to conserve the total angular momentum, the difference in the\neffective gravitational potential between the dust ($-GM\/r$) and the gas\n($-GM(1-2\\eta)\/r$), including the work done by the pressure gradient,\ncauses a net energy liberation. Part of the liberated energy is converted\ndirectly into thermal energy, and part is used for supplying energy to\nthe turbulence (see discussion in \\S \\ref{sec:ene-into-tur}).\n\nDust particle accretion occurs either individually or\ncollectively. Individual dust particles suffer gas drag given a\nvelocity difference between the particle and the gas. Usually the orbital\nvelocity of the gas is slower than that of the dust particle, and the gas drag\nforce decelerates the particle orbital motion. Thus, the particle loses its\nangular momentum and drifts inward. In addition to this individual drift,\na collective drag works on the whole dust layer. The dust layer\nat the midplane orbits faster than the gas layer at higher altitudes because in the\ndust layer the increased inertia of the enriched dust particles weakens the\neffect of the gas pressure. The dust layer rotates with a velocity close to\nthe Keplerian value. The upper gas layer, devoid of the inertia of the dust,\norbits more slowly than the midplane dust layer. If the gas exhibits turbulent\nviscosity, then the slower orbiting gas layer exerts a drag force on the dust\nlayer, and consequently the dust layer loses its angular momentum and accretes\ninward. We consider the energy release rate of the dust caused by \nindividual and collective drag.\n\n\\subsection{Dust Accretion Caused by Individual Drag}\n\\label{sec:inddrag}\n\nWe calculate the gravitational energy released when dust\nparticles accrete towards the star due to drag on individual particles. In this\ncalculation, we assume that the gas disk is in a laminar flow state, and that\ndust particles drift inward steadily. Any random motion of particles caused by\ngas turbulence is neglected. We expect that in turbulent gas\ndisks, the average motion of the dust particles can be estimated from the motion in\nthe laminar disk (Bai \\& Stone 2010). Because the individual drag is more\neffective than the collective drag when the midplane dust-to-gas ratio is less\nthan unity (see Fig. \\ref{fig:Sigma} below), we take the $f_{\\mathrm{mid}\n\\ll1$ limit in the following calculation. The calculation for a general $f_{\\mathrm{mid}}$ is described in Appendix \\ref{sec:app-ind}.\n\nIn a laminar disk, the particle drift velocity $v_{d,r}$ is given by equation\n(2.11) of NSH86. For $\\rho_{d}\\ll\\rho_{g}$, the particle radial velocity\nis (see also eq. [23] of Takeuchi \\& Lin (2002); note the factor of 2\ndifference in the definition of $\\eta$\n\\begin{equation}\nv_{d,r}=-\\frac{2{T_{s}}}{{T_{s}^{2}+1}}\\eta v_{\\mathrm{K}}\\ . \n\\label{eq:vdr\n\\end{equation}\nThe gas drifts in the opposite direction with a velocity $v_{g,r}$. From\nangular momentum conservation, the gas drift velocity is\n\\begin{equation}\nv_{g,r}=-\\frac{{\\rho_{d}}}{{\\rho_{g}}}v_{d,r}\\ . \n\\label{eq:vgr\n\\end{equation}\nThe effective gravities (including the pressure gradient force) acting on the\ndust and on the gas are $g_{d}=-GM\/r^{2}=-r\\Omega_{\\mathrm{K}}^{2}$ and\n$g_{g}=-GM(1-2\\eta)\/r^{2}=-r\\Omega_{g}^{2}$, respectively. While the total\nangular momentum is conserved, the total energy is not conserved because of the\ndifference in the effective gravities acting on the gas and on the dust,\n$\\rho_{d}{g_{d}v_{d,r}+\\rho}_{g}g_{g}v_{g,r}{\\neq0}$, providing a source of\nenergy for turbulence.\n\nNote that, in the above calculation, $v_{d,r}$ and $v_{g,r}$ are the\n\\textquotedblleft terminal velocities\\textquotedblright. This means that the\ngravitational accelerations of the dust and gas are balanced by the drag\nforces, and liberated energy translates directly into thermal energy. Thus,\nno energy would be contributed to turbulence. In reality, however, if the gas\ndisk is turbulent, a steady terminal velocity is not expected, and acceleration phases of the dust must occur, as\ndiscussed in \\S \\ref{sec:ene-into-tur}. In an acceleration phase, the work\ndone by gravity first provides kinetic energy, which can be translated\ninto turbulent energy. We estimate this energy input to turbulence. In the\nfollowing calculation, the energy liberated from the accreting dust is\nestimated using the terminal velocity described above for simplicity. Only a\ncertain fraction of the liberated energy goes into turbulence. Thus, the\nenergy input to turbulence is a factor $C_{\\mathrm{eff}}(<1)$ times the\nfollowing estimation. The factor $C_{\\mathrm{eff}}$ will be\ndetermined in \\S \\ref{sec:comp-sim} to be $\\approx0.2$ by comparing our\nmodel with the numerical simulation of turbulence in the dust layer\ndeveloped by JHK06.\n\nGiven that some portion of the accretion energy of the dust is consumed by the outward motion of the gas, the liberated gravitational energy per unit surface area of the disk is\n\\begin{equation}\n\\frac{\\partial E_{\\mathrm{drag}}}{\\partial t}=\\frac{1}{2}\\int_{-\\infty\n}^{\\infty}(\\rho_{d}{g_{d}v_{d,r}+}\\rho_{g}{g_{g}v_{g,r})dz}~,\n\\end{equation}\nwhere the factor of $1\/2$ comes from the fact that half of the work done by\ngravity is used for acceleration (and deceleration) of the azimuthal velocity\nof the dust (and of the gas) as their semi-major axes change. Using equation\n(\\ref{eq:vdr}) and (\\ref{eq:vgr}), the energy liberation rate reduces to\n\\begin{equation}\n\\frac{\\partial E_{\\mathrm{drag}}}{\\partial t}=2\\eta^{2}v_{\\mathrm{K}\n^{2}\\Omega_{\\mathrm{K}}T_{s}\\Sigma_{d,\\mathrm{drag}}\\ , \n\\label{eq:ene_drag\n\\end{equation}\nwhere the effective ``surface density'' $\\Sigma_{d,\\mathrm{drag}}$ of the dust\ni\n\\begin{equation}\n\\Sigma_{d,\\mathrm{drag}}=\\frac{1}{{T_{s}^{2}+1}}\\Sigma_{d}~.\n\\label{eq:sig_drag2\n\\end{equation}\nIn the above calculation, we assume $f_{\\mathrm{mid}}\\ll1$. The calculation of\nthe energy liberation rate for general $f_{\\mathrm{mid}}$ is described in\nAppendix \\ref{sec:app-ind}. In the numerical calculations in the subsequent\nsections, we use equation (\\ref{eq:sig_drag_ap}) for the effective surface\ndensity\n\\footnote{Equation (\\ref{eq:sig_drag}) does not coincide with\nequation (\\ref{eq:sig_drag2}) in the limit of $f_{\\mathrm{mid}}\\ll1$. This\ndiscrepancy comes from the fact that the calculation of equation\n(\\ref{eq:sig_drag}) includes correction terms of\nthe order of $\\eta^{2}$ in the drift velocity, while equation\n(\\ref{eq:sig_drag2}) considers only the terms of $\\eta$. The difference\nbetween equations (\\ref{eq:sig_drag2}) and (\\ref{eq:sig_drag}) is at most\nfactor 2 (for $T_{s}\\gg1)$, and is not significant.},\n\\begin{equation}\n\\Sigma_{d,\\mathrm{drag}}=\\frac{{\\Sigma_{d}}}{\\sqrt{\\pi}}\\int_{-\\infty\n^{\\infty}{\\frac{{\\exp(-\\tilde{z}^{2})}}{{T_{s}^{2}+\\beta^{2}}}\\left[\n1-\\frac{T_{s}^{2}}{2(T_{s}^{2}+\\beta^{2})}\\right] d\\tilde{z}}\\ .\n\\label{eq:sig_drag\n\\end{equation}\nIf the dust particles are tightly coupled to the gas ($T_{s}\\ll1$) and\ndust sedimentation is weak ($f_{\\mathrm{mid}}\\ll1$), then\n$\\Sigma_{d,\\mathrm{drag}}$ is simply equal to the dust surface density\n$\\Sigma_{d}$.\n\n\\subsection{Dust Accretion Caused by Collective Drag}\n\\label{sec:coldrag}\n\nNext, we consider dust accretion due to collective drag acting on the entire\ndust layer. Because the faster-orbiting dust particles drag the gas in the\ndust layer, the orbital velocity of the gas is largest at the midplane and\ndecreases with altitude. If the gas in the dust layer is turbulent, the\nvariation in orbital velocity $v_{g,\\theta}$ with altitude $z$ induces\nReynolds stress $P_{\\theta z}$. This causes a transfer of angular momentum from the dust\nlayer to the upper gas layer, resulting in accretion of the dust layer.\n\nTo calculate the energy liberation rate, we make a few key\nassumptions. First, we consider only the $\\theta z$-component of the Reynolds\nstress, $P_{\\theta z}$, neglecting the other components $P_{r\\theta}$ and\n$P_{zr}$. Ignoring $P_{r\\theta}$ means that turbulence in the dust layer does\nnot transfer the angular momentum efficiently in the radial direction. This is\nexpected for turbulence induced by hydrodynamical instabilities such as\nconvective instability (e.g., Stone \\& Balbus 1996; Lesur \\& Ogilvie 2010).\nThe Reynolds stress $P_{zr}$ is also neglected for simplicity. Brauer et al.\n(2007) pointed out that $P_{zr}$ changes the velocity profiles of the dust and\ngas from those derived by NSH86 by a factor $\\sim3$, and thus $P_{zr}$\ncannot be neglected in a rigorous discussion.\nObtaining the exact velocity profiles including $P_{zr}$ requires\nnumerical calculations. In this paper, we use the velocity profiles calculated\nanalytically by NSH86 for simplicity. This induces an error of a factor\n$\\sim3$ in our estimation of turbulence strength. Thus, our estimate is\nlimited to an order-of-magnitude argument.\n\nThe second assumption is that the turbulent layer has a thickness\ncomparable to that of the dust layer. If the turbulent layer were much thicker\nthan the dust layer and most of the volume of the turbulent layer were free of\nthe dust, then its structure would be controlled by the gas, \nunaffected by the properties of the dust and the structure of the dust layer.\nThe dust layer would behave just like a boundary wall at the bottom of\nthe turbulent layer. Turbulence in such a thick boundary gas layer has\nbeen discussed using the analogy of the Ekman layer (e.g., Cuzzi et al. 1993). \nIf the turbulent layer were dominated by the\ndust layer, then the structure of the dust layer would control the\nturbulence strength. We focus on such a dusty turbulent layer. Following\nYoudin \\& Chiang (2004), we consider the conditions needed for the\nturbulent and dust layers to be of similar thickness. The thickness of\nthe turbulent layer, from dimensional analysis, is the Ekman length,\n\\begin{equation}\nh_{E}\\sim\\sqrt{\\frac{\\nu}{\\Omega_{\\mathrm{K}}}}~, \n\\label{eq:he\n\\end{equation}\nwhere $\\nu$ is the turbulent viscosity, provided that the viscosity and the\nCoriolis force determine the layer structure. The thickness\nof the dust layer (eq. [\\ref{eq:hd}]), which is determined by the balance\nbetween sedimentation and diffusion of the dust particles, is\n\\begin{equation}\nh_{d}\\sim\\sqrt{\\frac{\\nu}{T_{s}\\Omega_{\\mathrm{K}}}}~, \n\\label{eq:hd2\n\\end{equation}\nwhere $\\nu \\sim D_g$ is used.\nFor small particles ($T_{s}\\la1$), $h_{d}$ is larger than $h_{E}$, meaning\nthat such small particles move further out of the turbulent layer and\nmodify the structure of the turbulent layer. It is expected that the thickness\nof the turbulent layer is not given by $h_{E}$, but by $h_{d}$ (see also S98;\nGoodman \\& Pindor 2000). For large particles ($T_{s}\\gg1$), the turbulent\nlayer is much thicker than the dust layer, and its thickness is expected to be\n$h_{E}$. In the following discussion, we consider a dusty turbulent\nlayer with a thickness similar to $h_{d}$. Thus, the analysis in this\npaper is probably not appropriate for large particles ($T_{s}\\gg1$).\n\nUnder the above assumptions, the energy liberation rate is estimated. Because\nthe collective drag is effective only if the midplane dust-to-gas ratio is\nlarger than unity (see Fig. \\ref{fig:Sigma} below), we consider the case in\nwhich a dense dust layer has formed ($f_{\\mathrm{mid}}\\ga1$) and use the\nplate drag approximation (Goldreich \\& Ward 1973; Goodman \\& Pindor 2000;\nWeidenschilling 2003). The calculation for general $f_{\\mathrm{mid}}$ is\ndescribed in Appendix \\ref{sec:app-col}. The Reynolds stress $P_{\\theta z}$\nnear the boundary between the dust layer (or the turbulent layer) and the gas\nlayer is estimated as\n\\begin{equation}\nP_{\\theta z}=\\rho_{g}\\nu\\frac{{\\partial v_{g,\\theta}}}{{\\partial z}}\\sim\n-\\rho_{g}\\nu\\frac{\\eta v_{\\mathrm{K}}}{h_{d}}~. \n\\label{eq:Ptz\n\\end{equation}\nThis stress extracts angular momentum from the dust layer and transfers it to\nthe gas layer. The unit surface of the dust layer loses angular momentum $\\partial\nL_{d}\/\\partial t=rP_{\\theta z}$, and the corresponding energy change is $\\partial\nE_{d}\/\\partial t=\\Omega_{\\mathrm{K}}\\partial L_{d}\/\\partial t$. The gas layer\ngains the same amount of angular momentum $\\partial L_{g}\/\\partial\nt=-rP_{\\theta z}$, but the energy change is different from that of the\ndust layer because of the work done by the pressure gradient: $\\partial\nE_{g}\/\\partial t=\\Omega_{g}\\partial L_{g}\/\\partial\nt=(1-\\eta)\\Omega_{\\mathrm{K}}\\partial\nL_{g}\/\\partial t$. In total, the energy liberation rate (the minus sign is\nadded), using equations (\\ref{eq:rhog}),(\\ref{eq:hd}),(\\ref{eq:fdmid}), and\n(\\ref{eq:Ptz}), is\n\\begin{equation}\n\\frac{\\partial E_{\\mathrm{vis}}}{\\partial t}=-(\\Omega_{\\mathrm{K}}-\\Omega\n_{g})rP_{\\theta z}=2\\eta^{2}v_{\\mathrm{K}}{}^{2}\\Omega_{\\mathrm{K}}T_{s\n\\Sigma_{d,\\mathrm{vis}}~, \n\\label{eq:ene_vis\n\\end{equation}\nwhere\n\\begin{equation}\n\\Sigma_{d,\\mathrm{vis}}=\\frac{1}{2\\sqrt{2\\pi}}\\frac{1+2T_{s}}{1+T_{s}}\\frac\n{1}{f_{\\mathrm{mid}}}\\Sigma_{d}~. \n\\label{eq:sig_vis-approx\n\\end{equation}\nIn the above calculation, the viscosity coefficient is modeled as\n$\\nu=\\alpha c_g h_g$. The same $\\alpha$ is used for both $D_g$ and\n$\\nu$ for simplicity. \nThe effective surface density $\\Sigma_{d,\\mathrm{vis}}$ is inversely\nproportional to the midplane dust-to-gas ratio $f_{\\mathrm{mid}}$ and only\nweakly depends on $T_{s}$. Note that the above relationship is derived for\n$f_{\\mathrm{mid}}\\ga1$. The effective surface density for general\n$f_{\\mathrm{mid}}$ is derived in Appendix \\ref{sec:app-col} and is given by\n\\begin{eqnarray}\n\\Sigma_{d,\\mathrm{vis}} &=&\n\\frac{C_{\\mathrm{str}}}{\\sqrt{{\\pi}}}f_{\\mathrm{mid}}\\Sigma_{d}\\frac{{1+2T}_{s}}{1+{T_{s}}}\n\\nonumber \\\\\n& \\times & \\int_{-\\infty}^{\\infty}{\\frac\n{{\\beta^{2}+2\\beta T_{s}^{2}-T_{s}^{2}}}{{(T_{s}^{2}+\\beta^{2})^{2}}\n\\frac{{\\tilde{z}^{2}\\exp(-2\\tilde{z}^{2})}}{{[1+C_{\\mathrm{str}\nf_{\\mathrm{mid}}\\exp(-\\tilde{z}^{2})]}}d\\tilde{z}}\\ , \\nonumber \\\\\n\\label{eq:sig_vis\n\\end{eqnarray}\nwhere\n\\begin{equation}\n{C_{\\mathrm{str}}=}\\frac{1}{T_{s}^{2}+1}~.\n\\end{equation}\nExpression (\\ref{eq:sig_vis}) is complicated, but for large $f_{\\mathrm{mid}}$, its dependence on $T_{s}$ and $f_{\\mathrm{mid}}$ is similar to that of\nthe simpler equation (\\ref{eq:sig_vis-approx}); $\\Sigma_{d,\\mathrm{vis}}\\propto\nf_{\\mathrm{mid}}^{-0.9}$ and depends only weakly on $T_{s}$.\n\n\\subsection{Energy Dissipation in Turbulence}\n\nThe turbulent energy of the largest eddies transfers to smaller eddies,\nand finally dissipates to thermal energy via decay of the smallest eddies due\nto molecular viscosity. The size and velocity of the largest eddies are\nassumed to be $l_{g,\\mathrm{eddy}}=\\alpha^{1\/2}h_{g}$ and $u_{g,\\mathrm{eddy\n}=\\alpha^{1\/2}c_{s}$. This assumption means that the turnover time of the\nlargest eddies is the Keplerian time ($\\tau_{g,\\mathrm{eddy}\n=l_{g,\\mathrm{eddy}}\/u_{g,\\mathrm{eddy}}=\\Omega_{\\mathrm{K}}^{-1}$; Cuzzi et\nal. 2001). Using the Kolmogorov scaling law, the energy dissipation rate per\nunit volume and unit time is\n\\begin{equation}\n\\frac{\\partial\\varepsilon_{\\mathrm{turb}}}{\\partial t}=({\\rho_{g\n+C_{\\mathrm{ene}}\\rho_{d})}\\frac{{u_{g,\\mathrm{eddy}}^{2}}}{\\tau\n{_{g,\\mathrm{eddy}}}}=({\\rho_{g}+C_{\\mathrm{ene}}\\rho_{d})}\\alpha h_{g\n^{2}\\Omega_{\\mathrm{K}}^{3}\\ ,\n\\end{equation}\nwhere the factor $C_{\\mathrm{ene}}$ represents the fact that dust\nparticles that are coupled only weakly to the gas do not contribute to\nthe turbulent energy, and is given by (see Appendix \\ref{sec:app-ene-dis})\n\\begin{equation}\nC_{\\mathrm{ene}}=\\left\\{\n\\begin{array}\n[c]{ccc\n\\frac{1}{T_{s}T_{e}+1} & \\mathrm{for} & T_{s}\\,\\leq T_{e}\\\\\n\\frac{1}{T_{s}^{2}(T_{e}^{-2}+1)} & \\mathrm{for} & T_{s}\\,>T_{e\n\\end{array}\n\\right. \\ ,\n\\end{equation}\nwhere $T_{e}=\\tau_{g,\\mathrm{eddy}}\\Omega_{\\mathrm{K}}$ is the non-dimensional\nturnover time of the largest eddies. For most of this paper, we\nconsider $T_{e}=1$. The effect of the dust inertia is ignored for simplicity.\n\nEnergy dissipation of turbulence occurs only in the dust layer. This\nassumption may be problematic for $T_{s}\\gg1$, as discussed in\n\\S \\ref{sec:coldrag}, but for simplicity, we assume that energy dissipation\noccurs in $-\\sqrt{2}h_{d}1$, $\\alpha$ peaks\nat a larger $Z_{\\mathrm{disk}}$ (or $f_{\\mathrm{mid}}$).\n\nFigure \\ref{fig:fdmid} shows the variation in the midplane dust-to-gas ratio\n$f_{\\mathrm{mid}}$ with the disk metallicity $Z_{\\mathrm{disk}}$. For small particles ($T_{s}\\la1$), $\\alpha$ is\nproportional to $T_{s}$, and the midplane dust-to-gas ratio\n$f_{\\mathrm{mid}}\\propto(T_{s}\/\\alpha)^{1\/2}$ does not depend on\n$T_{s}$. For a small dust-to-gas ratio ($f_{\\mathrm{mid}}\\ll1$),\n$f_{\\mathrm{mid}}\\propto Z_{\\mathrm{disk}}^{2\/3}$, and for $f_{\\mathrm{mid}}\\gg1$,\n$f_{\\mathrm{mid}}\\propto Z_{\\mathrm{disk}}^{1.9}$. Note that if the disk metallicity is larger than $0.2$, then the midplane dust-to-gas ratio\nbecomes as large as $100$, which is high enough to induce a gravitational\ninstability in the dust layer. This result is consistent with the previous\nworks by, e.g., S98. Figure \\ref{fig:hd} shows that the dust layer thickness $h_{d}$\nis of the order of $10^{-3}-10^{-2}$ times the thickness of the gas disk.\n\n\\begin{figure}[ptb]\n\\epsscale{1.1} \\plotone{f4.eps}\\caption{The scale height of the dust layer\n$h_{d}$ normalized by the gas scale height $h_g$. The solid lines are\n calculated for various values of the stopping time $T_{s}$ by our\n model. The dashed line is taken from the S98 model, which is\n calculated using equation (\\ref{eq:hd-sek}).}\n\\label{fig:hd\n\\end{figure}\n\n\\section{COMPARISONS WITH PREVIOUS STUDIES}\n\\label{sec:comppre}\n\nIn this section, the dust layer model that is based on the energetics of\ndust accretion is compared with results from previous work. Previous studies \nhave analyzed the detailed physics of the dust layer, including the\nonset of KH or streaming instabilities, and several stabilizing effects, such as\nKeplerian shear, using both analytical and numerical approaches. It is of\ninterest to determine which properties of the dust layer are reproduced by our\nmodel and which are missing.\n\n\\subsection{Comparison with Previous Analytical Studies}\n\n\\subsubsection{Comparison with Sekiya (1998)}\n\nS98 analytically solved the density structure of the dust layer, assuming that\nthe layer structure was adjusted to be marginally unstable to KH instability.\nTo obtain the density structure, S98 assumed that the Richardson number remained \nclose to the critical value for instability throughout the dust\nlayer. S98 considered small dust particles that were tightly coupled to the gas \n($T_{s}\\ll1$), and treated the gas and the dust as a single fluid. Then, the\ndensity structure of the dust layer was calculated for various values of the\ndisk metallicity $Z_{\\mathrm{disk}}$. In S98, the dust layer\nstructure was determined by the argument for the stability of a\nstratified fluid, not \nby the balance between sedimentation and diffusion of the dust. However, it is\npossible to interpret the result of S98 as follows: The turbulent\nstrength in the dust layer is adjusted such that the turbulent diffusion of\nthe dust maintains the marginally unstable density structure. Using this\ninterpretation, we calculate the effective value of the turbulent diffusion\nparameter $\\alpha$ as a function of $Z_{\\mathrm{disk}}$ for $T_{s}\\ll1$.\n\nFrom equation (22) of S98, the midplane dust-to-gas ratio,\n$f_{\\mathrm{mid}}$, is related to the disk metallicity,\n$Z_{\\mathrm{disk}}$, as\n\\begin{eqnarray}\n{Z_{\\mathrm{disk}}} &=& \\sqrt{\\frac{{2\\mathrm{{Ri}}}\\tilde{\\eta}}{\\pi}} \\left\\{\n\\ln\\left[ {\\left( {1+\\sqrt{1-{{\\left( {\\frac{1}{{1+{f_{\\mathrm{mid}}}}}\n}\\right) }^{2}}}}\\right) \\left( {1+{f_{\\mathrm{mid}}}}\\right) }\\right]\n\\right. \\nonumber \\\\\n& &\\left. -\\sqrt{1-{{\\left( {\\frac{1}{{1+{f_{\\mathrm{mid}}}}}}\\right) }^{2}}}\n\\right\\} \\ , \n\\label{eq:fd-sek\n\\end{eqnarray}\nwhere the Richardson number $\\mathrm{{Ri}}$ is constant\nthroughout the dust layer, and the self-gravity of the dust layer is neglected.\nThe dust density distribution is given by\n\\begin{equation}\n\\frac{{{\\rho_{d}}}}{{{\\rho_{g}}}}={\\left[ {\\frac{{{z^{2}}}}{{\\mathrm{{Ri\n}{\\tilde{\\eta}h}}_{g}^{2}}+{{\\left( {\\frac{1}{{1+{f_{\\mathrm{mid}}}}\n}\\right) }^{2}}}\\right] ^{-\\frac{1}{2}}}-1\\ . \n\\label{eq:dendist-sek\n\\end{equation}\nThe half-thickness of the dust layer, $z_{d}$, at which the dust density\nbecomes zero, is given by\n\\begin{equation}\nz_{d}=\\sqrt{\\mathrm{{Ri}}\\tilde{\\eta}\\left[ {1-(1+f_{\\mathrm{mid}})\n^{-2}\\right] }\\mathrm{\\ }h_{g}~. \n\\label{eq:zd\n\\end{equation}\nTo compare with the scale height of Gaussian distribution, $h_{d}$, we define\nthe scale height of the dust distribution as the vertical dispersion of the\ndust particles,\n\\begin{equation}\nh_{d,\\mathrm{{Sek}}}^{2}=\\frac{\\int_{0}^{z_{d}}z^{2}\\rho_{d}dz}{\\int\n_{0}^{z_{d}}\\rho_{d}dz}\\ . \n\\label{eq:hd-sek\n\\end{equation}\nThe turbulent parameter $\\alpha_{\\mathrm{{Sek}}}$ is estimated from equation\n(\\ref{eq:hd}), using $T_{s}\\ll1$, as\n\\begin{equation}\n\\alpha_{\\mathrm{{Sek}}}=\\left( \\frac{h_{d,\\mathrm{{Sek}}}}{h_{g}}\\right)\n^{2}T_{s}\\ . \n\\label{eq:alpha-sek\n\\end{equation}\nFor $f_{\\mathrm{mid}}\\ga1$, equation (\\ref{eq:zd}) reduces to $z_{d\n\\approx(\\mathrm{{Ri}}\\tilde{\\eta})^{1\/2}h_{g}$. Substituting this expression\ninto $h_{d,\\mathrm{Sek}}\\approx z_{d}$ of equation (\\ref{eq:alpha-sek}) gives\n\\begin{equation}\n\\alpha_{\\mathrm{{Sek}}}\\approx\\mathrm{{Ri}}{\\tilde{\\eta}T_{s}\n\\ \\ \\ \\mathrm{for}\\ \\ \\ f_{\\mathrm{mid}}\\ga1\\ , \n\\label{eq:alpha-sek2\n\\end{equation}\nwhich can be compared with equation (\\ref{eq:alpha2}). For $f_{\\mathrm{mid}}\\ll1$, equation (\\ref{eq:fd-sek}) is approximated in the lowest order of\n$f_{\\mathrm{mid}}$ as\n\\begin{equation}\n{Z_{\\mathrm{disk}}}=\\frac{4}{3}\\sqrt{\\frac{{\\mathrm{{Ri}}}\\tilde{\\eta}}{\\pi\n}\\left[ {f_{\\mathrm{mid}}^{3\/2}+O(f_{\\mathrm{mid}}^{5\/2})}\\right] .\n\\end{equation}\nThus, the midplane dust-to-gas ratio is\n\\begin{equation}\nf_{\\mathrm{mid}}=\\left( { \\frac{9}{16}\\frac\n{\\pi}{{\\mathrm{{Ri}}\\tilde{\\eta}}}Z_{\\mathrm{disk}}^{2}}\\right)^{\\frac\n{1}{3}}\\ . \n\\label{eq:fdmid-sek\n\\end{equation}\nFor $f_{\\mathrm{mid}}\\ll1$ and $z \\ll h_g$, the dust density distribution\n(eq.[\\ref{eq:dendist-sek}]) is\n\\begin{equation}\n\\frac{\\rho_{d}}{\\rho_{g}}\n \\approx{f_{\\mathrm{mid}}}\\left[ {1-\\frac{{1}}{{2\\mathrm{{Ri}\n{\\tilde{\\eta}f_{\\mathrm{mid}}}}}\\left( \\frac{z}{h_{g}}\\right) }^{2\n \\right] \\ .\n\\end{equation}\nTo the order of $(z\/h_{d})^{2}$, this distribution is approximated by the\nGaussian distribution $f_{\\mathrm{mid}}\n\\exp[-z^{2}\/(2h_{d,\\mathrm{Sek}}^{2})] \\approx f_{\\mathrm{mid}}\n [1-z^{2}\/(2h_{d,\\mathrm{Sek}}^{2})]$, and its scale height $h_{d,\\mathrm{Sek}}$ is \n\\begin{equation}\nh_{d,\\mathrm{Sek}}\\approx\\sqrt{\\mathrm{{Ri}}\\tilde{\\eta}{f_{\\mathrm{mid}}\n}h_{g}\\approx{\\left[ {\\frac{{\\mathrm{{3Ri}}}}{4}{\\pi^{1\/2}}{Z_{\\mathrm{disk}}}\\tilde{\\eta}}\\right] ^{\\frac{1}{3}}}{h_{g}}\\ .\n\\end{equation}\nThe turbulent parameter $\\alpha_{\\mathrm{Sek}}$ of equation\n(\\ref{eq:alpha-sek}) is\n\\begin{equation}\n\\alpha_{\\mathrm{Sek}}\\approx{\\left[ {\\frac{{\\mathrm{{3Ri}}}}{4}{\\pi^{1\/2\n}{Z_{\\mathrm{disk}}}\\tilde{\\eta}}\\right] ^{\\frac{2}{3}}}{T_{s\n}\\ \\ \\ \\ \\mathrm{for}\\ \\ \\ f_{\\mathrm{mid}}\\ll1\\ . \n\\label{eq:alpha-sek3\n\\end{equation}\nComparing this expression to equation (\\ref{eq:alpha1}) provides a relationship\nbetween the energy supply efficiency to turbulence, $C_{\\mathrm{eff}}$, and\nthe critical Richardson number, $\\mathrm{Ri}$, as ${C_{\\mathrm{eff}\n}=(3\/4)\\mathrm{{Ri}}$. In the above discussion on $\\alpha_{\\mathrm{Sek}}$, we\nused the dust layer thickness $h_{d,\\mathrm{Sek}}$ derived from the Gaussian\nfit. However, the thickness defined by the vertical dispersion\n(eq.[\\ref{eq:hd-sek}]) in the small $f_{\\mathrm{mid}}$ limit is $0.63$ times\nthinner than that defined by the Gaussian fit. Thus, the energy supply\nefficiency is $0.4=0.63^{2}$ times smaller than the above estimate and thus\n$C_{\\mathrm{eff}}\\approx0.3\\mathrm{{Ri}}$. The relationship between\n$C_{\\mathrm{eff}}$ and $\\mathrm{Ri}$ can also be derived by comparing the\nexpressions for $f_{\\mathrm{mid}}$ (eqs. [\\ref{eq:fdmid1}] and\n[\\ref{eq:fdmid-sek}]) as ${C_{\\mathrm{eff}}}=16\/(9\\pi^{3\/2})\\mathrm{{Ri\n}\\approx0.3\\mathrm{{Ri}}$. We compare our result with the numerical simulation\ndeveloped by JHK06, which will be discussed in \\S \\ref{sec:comp-sim}. The comparison\nwith this simulation suggests a slightly smaller energy supply efficiency,\n\\begin{equation}\nC_{\\mathrm{eff}}\\approx0.24\\mathrm{{Ri}}\\ . \n\\label{eq:ceff-ri\n\\end{equation}\n\n\nThe turbulent parameter $\\alpha_{\\mathrm{Sek}}$ given by equation\n(\\ref{eq:alpha-sek}) is numerically calculated and is plotted using dashed lines\nin Figure \\ref{fig:alpha}. The critical Richardson number $\\mathrm{Ri}=0.8$ is\nadopted to fit the simulation results provided by JHK06. For $Z_{\\mathrm{disk}}\\ll1$\nand $T_{s}\\la1$, the solid and dashed lines agree with each other very\nwell. This is expected from the above discussion that the dependence of\n$\\alpha$ and $\\alpha_{\\mathrm{Sek}}$ on $Z_{\\mathrm{disk}}$ and $T_{s}$ for\nsmall $f_{\\mathrm{mid}}$ are the same (eqs. [\\ref{eq:alpha1}] and\n[\\ref{eq:alpha-sek3}]). For $Z_{\\mathrm{disk}}\\ga 10^{-2}$, a deviation\nbetween the solid and dashed lines appears, and the difference increases\nwith $Z_{\\mathrm{disk}}$. We note that for $Z_{\\mathrm{disk}} \\ga 0.1$,\nthe S98 density distribution shows a cusp at the midplane, which is\na significant difference from the Gaussian distribution of our model. In\nFigure \\ref{fig:fdmid}, the midplane dust-to-gas ratio of the S98 model\nis plotted with a dashed line. For small $Z_{\\mathrm{disk}}\\la0.05$, the\nresult from S98 is consistent with our result (the dashed line coincides\nwith the solid lines for $T_{s}\\la1$). For \n$Z_{\\mathrm{disk}}\\ga0.05$, the midplane dust-to-gas ratio from the S98 model\nis larger than our results. Both results are still qualitatively consistent,\nindicating that, for $Z_{\\mathrm{disk}}\\ga0.1$, the midplane dust-to-gas\nratio becomes large enough for gravitational instability ($f_{\\mathrm{mid}}\\sim100$).\n\n\\subsubsection{Comparison with Michikoshi \\& Inutsuka (2006)}\n\n\\begin{figure}[ptb]\n\\epsscale{1.1} \\plotone{f5.eps}\\caption{Comparison with the results in\nMichikoshi \\& Inutsuka (2006). The solid line shows the relationship between the\nmidplane dust-to-gas ratio $f_{\\mathrm{mid}}$ and the particle radius $a$.\nThe dashed line shows the boundary at which the growth rate of the KH instability\nis $\\omega=0.1$ (above the dashed line, $\\omega>0.1$). We use the same disk\nmodel that is adopted in Michikoshi \\& Inutsuka (2006). The dashed line is read from\nFig. 16 in Michikoshi \\& Inutsuka (2006).\n\\label{fig:michi\n\\end{figure}\n\nMichikoshi \\& Inutsuka (2006) analyzed the growth rate of the KH instability\nof a dust layer, taking into account the relative motion and the friction\nbetween the dust and the gas. Their formulation does not assume tight dust-gas\ncoupling, and thus it can be applied to scenarios involving large dust particles ($T_{s}\\gg1$),\nwhile the vertical gravity, Coriolis force, and Keplerian shear are neglected.\nThe initial velocity gradient in the vertical direction is caused by the dust\ninertia, using the formula provided in NSH86. The growth rate of instability has been derived\nfor a wide range of dust sizes $a$ and midplane dust-to-gas\nratios $f_{\\mathrm{mid}}$, and is summarized in their Figure 16. They argued\nthat if the effect of Keplerian shear is taken into account, the line\ncorresponding to the growth rate $\\omega=1$ in the $a$-$f_{\\mathrm{mid}}$\nplane would be the boundary between the stable and unstable\nconfigurations. This marginally unstable structure of the dust layer\npresented by Michikoshi \\& Inutsuka (2006) can be considered an\nextension of the S98 model for general sizes of the dust particles, and\nprovides a reference to be compared with our model.\n\nIn Figure \\ref{fig:michi}, we compare our model in the $a$-$f_{\\mathrm{mid}}$\nplane with the line at which the growth rate $\\omega$ has a constant critical\nvalue. We adopt the critical growth rate as $\\omega=0.1$, which is smaller\nthan the value Michikoshi \\& Inutsuka (2006) suggested. The qualitative\nbehavior does not differ between $\\omega=0.1$ and $\\omega=1$. Our model well\nreproduces the result of Michikoshi \\& Inutsuka (2006). The agreement of our model\nwith the results of S98 and of Michikoshi \\& Inutsuka (2006) suggests that the\nonset of KH instability is controlled by the energy supply due to dust\naccretion over a wide range of particle sizes, $a$, or stopping times,\n$T_{s}$. It must be noted, however, that the initial state assumed in Michikoshi \\&\nInutsuka (2006) is such that the vertical velocity shear appears only\nin the dust layer, neglecting the velocity shear in the Ekman-like boundary\nlayer that may appear between the dust layer and the upper gas layer. This is\nthe same assumption that we adopt. As discussed in \\S \\ref{sec:coldrag}, this\nassumption is appropriate for small dust particles ($T_{s}\\la1$), while for\nlarge particles ($T_{s}\\gg1$) the vertical velocity profile of the thick turbulent \nboundary layer may quickly deviate from that given by NSH86. Though the\nagreement with the result by Michikoshi \\& Inutsuka (2006) shows the applicability\nof our model to KH instability in the dust layer for general values of $T_{s\n$, we should keep in mind the limitation of the models mentioned above.\n\n\\begin{figure}[ptb]\n\\epsscale{1.1} \\plotone{f6.eps}\\caption{The turbulent viscosity parameter $\\alpha$\nestimated from our model, compared with the simulation described in JHK06. The solid\nline is calculated from our model and the dashed line is taken from the S98 model. The\nsquares are the simulation results, $\\delta_{t}$ in Table 2 of JHK06. The\nstopping time $T_{s}=0.1$ is adopted both for our model and the simulation.\n\\label{fig:alpha-joh\n\\end{figure}\n\n\\subsection{Comparison with Previous Numerical Simulations}\n\\label{sec:comp-sim}\n\n\\subsubsection{Comparison with Johansen et al. (2006)}\n\nNumerical simulations of KH instability induced by formation of\na dust layer were presented in JHK06. Their two-dimensional \nsimulation on the $\\theta z$-plane solved for\nboth the gas and dust motions. They obtained a quasi-steady or oscillating\ndensity distribution of the dust in which dust settling and turbulent\ndiffusion balanced each other. The diffusion coefficient in the vertical\ndirection, $\\delta_{t}$, was measured from the scale height of the dust layer\nand was summarized in their Table 2. In Figure \\ref{fig:alpha-joh}, the\nturbulent viscosity parameter $\\alpha$ of equation (\\ref{eq:alpha}) (the solid\nline) is compared with the diffusion coefficient measured from the simulation\n(the squares). To fit the simulation result, the efficiency parameter of the\nenergy supply to turbulence, $C_{\\mathrm{eff}}$, is set to $0.19$. We also\nfit the turbulent diffusion parameter for the density distribution of S98,\n$\\alpha_{\\mathrm{Sek}}$, to the simulation. We adopt the critical Richardson\nnumber as $\\mathrm{Ri}=0.8$ for fitting. This is slightly smaller than the\nvalue JHK06 suggested ($\\mathrm{Ri}=1.0$). It is seen that the numerical\nsimulation is well explained both by our model and the S98 model. This means that the\ndust layer is maintained so as to provide a constant Richardson number, and in\nthe parameter range that JHK06 surveyed ($0.01 h_{d}$).\nFrom equation (\\ref{eq:hd}) and $l_{g,\\mathrm{eddy}}=\\alpha^{1\/2}h_{g}$, the\ncondition for the dust layer thickness to be larger than the largest eddy size\n($h_{d}>l_{g,\\mathrm{eddy}}$) is $T_{s}<1\/\\sqrt{2}$. Thus, for $T_{s\n>1\/\\sqrt{2}$, our model results in a dust layer that is too thick (or an eddy\nsize that is too small). One possible remedy for this situation is to remove the assumption that\nthe turnover time of the largest eddies is equal to the Keplerian time. By\nintroducing a parameter $\\xi$ ($0\\leq \\xi \\leq1\/2$), the largest eddy size and the\nvelocity are expressed as $l_{g,\\mathrm{eddy}}=\\alpha^{(1\/2)+\\xi}h_{g}$ and\n$u_{g,\\mathrm{eddy}}=\\alpha^{(1\/2)-\\xi}c_{s}$. The non-dimensional turnover time\nof the largest eddies is $T_{e}=(l_{g,\\mathrm{eddy}}\/u_{g,\\mathrm{eddy\n})\\Omega_{\\mathrm{K}}=\\alpha^{2\\xi}$. The dust layer thickness for $T_{e}\\neq1$\nis provided by equation (21) of YL07,\n\\begin{equation}\n{h_{d}}=\\sqrt{\\frac{\\alpha}{{{T_{s}}}}}{\\left( {1+\\frac{{{T_{s}}{T_{e}}^{2}\n}{{{T_{s}}+{T_{e}}}}}\\right) ^{-1\/2}}{h_{g}}\\ .\n\\end{equation}\nFor $T_{s}>1\/2$, we impose the condition that the largest eddy size\n$l_{g,\\mathrm{eddy}}=\\alpha^{(1\/2)+\\xi}h_{g}$ is equal to the above $h_{d}$.\nThis condition determines the turnover time by the equation\n\\begin{equation}\n{T_{s}}{T_{e}}({T_{s}}+{T_{e}}+{T_{s}}{T_{e}}^{2})-{T_{s}}-{T_{e}}=0\\ .\n\\label{eq:Teddy\n\\end{equation}\nUsing this $T_{e}$ (or $\\xi$), the turbulent viscosity parameter $\\alpha$ and\nthe dust layer thickness $h_{d}$ are recalculated and plotted with the dashed\nline in Figure \\ref{fig:hd-ts-joh}. As expected, the introduction of a new\nparameter ${T_{e}}$ suppresses the largest eddy size and the turbulent\ndiffusion of dust particles for $T_{s}>1\/\\sqrt{2}$, improving the comparison with the\nsimulation result. However, we need to determine whether the turnover\ntime of the largest eddies, $T_{e}\\Omega_{\\mathrm{K}}^{-1}=\\alpha^{2\\xi\n\\Omega_{\\mathrm{K}}^{-1}$, in the simulation presented in JHK06 for $T_{s}=1$ is\nactually smaller than the Keplerian time, as equation (\\ref{eq:Teddy})\npredicts. Cuzzi et al. (1993) argue that, based on laboratory\nmeasurements, the eddy turnover time in an Ekman layer would be\nsmaller than the Keplerian time by a factor $20-80$. The turbulent\nlayer possibly behaves as an Ekman layer for $T_s >1$, as discussed in\n\\S\\ref{sec:coldrag}. We are currently performing numerical\nsimulations using the same conditions as JHK06 to investigate the\nturbulence for $T_{s} \\ga 1$ in more detail (Ishitsu et al., in preparation).\n\nLee et al. (2010) performed a three-dimensional numerical simulation of the\nonset of KH instability. They solved simplified equations in which the\ndust and the gas were treated as a single fluid, but they included the effect of\nthe Keplerian shear in the radial direction. They found that the radial shear\nstabilizes the KH instability, and the critical Richardson number for instability\nis not always the standard value, $0.25$, but can be much smaller if the\nstabilizing effect of the radial shear is significant (see also Ishitsu \\&\nSekiya 2003). Equation (32) in Lee et al. (2010) shows that the ratio of the\nstabilizing effect by the radial shear to the destabilizing effect by the\nvertical shear is proportional to $\\mathrm{Ri}(1+f_{\\mathrm{mid}})\/f_{\\mathrm{mid}}\\approx\\mathrm{Ri}\/f_{\\mathrm{mid}}$ for\n$f_{\\mathrm{mid}}\\ll1$, and thus the critical Richardson number should scale\nas $\\mathrm{Ri}\\propto f_{\\mathrm{mid}}$. In our model, we assume that the\nenergy supply efficiency $C_{\\mathrm{eff}}$, which is proportional to the\ncritical Richardson number (eq.[\\ref{eq:ceff-ri}]), is constant. However, the\nsimulation in Lee et al. (2010) suggests that\n$C_{\\mathrm{eff}}$ should also be proportional to $f_{\\mathrm{mid}}$\n(even for $f_{\\mathrm{mid}} \\ga 1$). If this is the case,\ndependence of $f_{\\rm mid}$ on $Z_{\\rm disk}$ would be milder than\nthat shown in Figure \\ref{fig:fdmid} ($f_{\\rm mid} \\propto Z_{\\rm\ndisk}^{1\/2}$ for $f_{\\rm mid} \\ll 1$ and $f_{\\rm mid} \\propto Z_{\\rm\ndisk}^{0.97}$ for $f_{\\rm mid} \\gg 1$).\nThe effect of the radial shear was not included\nin the simulation presented in JHK06, with which we compared our model\nin detail, because their simulation was two-dimensional in the $\\theta\nz$-plane. Extending JHK06 to three dimensions and including the radial\nshear effect are crucial to determining how the energy supply efficiency\n$C_{\\mathrm{eff}}$ behaves as $f_{\\mathrm{mid}}$ varies.\n\n\\subsubsection{Comparison with Bai \\& Stone (2010)}\n\\label{sec:BS10}\n\n\\begin{figure}[ptb]\n\\epsscale{1.1} \\plotone{f9a.eps} \\plotone{f9b.eps} \n\\caption{The turbulent\nviscosity parameter $\\alpha$ estimated from our model for $T_{s}=10^{-1}$ and 1\nis compared with the simulation presented in Bai \\& Stone (2010). The solid line is\ncalculated assuming the turnover time of the largest eddies is the Keplerian\ntime ($T_{e}=1$). The dashed line includes the variable $T_{e}$ calculated from\nequation (\\ref{eq:Teddy}). $(a)$ Comparison with the vertical diffusion\ncoefficient, $D_{g,z}(3\\mathrm{D})$, measured in the simulation (Table 2 of\nBai \\& Stone 2010). The squares and circles are the simulation results, \nconnected by lines that indicate which runs have the same size distribution\nof dust particles. The label of run R$xy$ means the particle size\ndistribution is such that the stopping time ranges $T_{s}\\in\\lbrack\n10^{-x},10^{-y}]$, e.g., in run R30, $T_{s}\\in\\lbrack10^{-3},10^{0}]$. In each\nrun, the disk metallicity $Z$ is $0.01$ and $0.03$. To compare the results with our model,\nwe assume that only the largest particles contribute to turbulence, and that the\neffective metallicity is estimated by $Z_{\\mathrm{disk}\n=Z\/N_{\\mathrm{type}}$, where $N_{\\mathrm{type}}=3$ (R10 and R21) or 7 (R30 and\nR41) is the number of particle species. We compare our model of $T_{s}=1$ with\nruns R10 and R30 (plotted in blue), and the model of\n$T_{s}=10^{-1}$ with runs R41 and R21 (plotted in red). $(b)$\nComparison with the radial diffusion coefficient, $D_{x}$, measured in the 3D\nruns of the simulation. We assume that the diffusion coefficient of the gas is\nrepresented by that of the smallest particles in the simulation. The values of\n$D_{x}$ for the smallest particles are read from Fig.9 of Bai \\& Stone (2010).\n\n\\label{fig:alpha-bai\n\\end{figure}\n\nBai \\& Stone (2010) performed a three-dimensional simulation, focusing on\ninvestigating turbulence induced by streaming instability. They found that\nthe streaming instability induced turbulence before the KH instability set in.\nThey measured the turbulent diffusion coefficient. However, it is difficult to\ncompare our model directly with the results in Bai \\& Stone (2010), because their\nsimulation includes particles of several sizes (3-7 species) while our model\nconsiders only single-sized particles. In the simulation, it was reported that\nonly large particles were responsible for inducing turbulence. Our model also\nshows that the energy liberation per unit mass of the dust is higher for\nlarger particles (it is proportional to $T_{s}$ for $T_{s}\\la1$). In order to\ncompare results, we assume that in the simulation, the turbulence is induced\nonly by the largest particles (i.e., particles of largest $T_{s}$). For\nexample, in the R41 run (in which the stopping time of the particles ranges\nfrom $10^{-4}$ to $10^{-1}$), we assume that only $T_{s}=0.1$ particles\nare responsible for turbulence. We then compare the simulation result with our\nmodel of $T_{s}=0.1$. In the simulation, each species has the same amount of\nmass. Because we consider the largest particles only, the total amount of the\ndust participating in driving turbulence is $\\Sigma\n_{d}\/N_{\\mathrm{type}}$, where $N_{\\mathrm{type}}$ is the number of particle\nspecies in the simulation. For example, in the R41 run, $N_{\\mathrm{type}}=7$,\nand we compare the simulation with a disk metallicity $Z=0.01$ with our\nmodel of $Z_{\\mathrm{disk}}=Z\/N_{\\mathrm{type}}=1.43\\times10^{-3}$. \n\nFigure \\ref{fig:alpha-bai}$a$ shows a comparison of the diffusion\ncoefficient obtained from the simulation ($D_{g,z}(\\mathrm{3D})$ in\nTable 2 of Bai \\& Stone (2010)) with our model (calculated with the\nparameter $\\tilde{\\eta}=0.05^{2}$, which was adopted by Bai \\& Stone\n(2010)). Although our model of $T_{s}=0.1$ agrees with simulations R41\nand R21 (plotted in red), we note a qualitative discrepancy between the $T_{s}=1$ model\nand simulations R30 and R10 (plotted in blue). The simulations indicate that the diffusion\ncoefficient decreases with the disk metallicity, and that its value for\n$Z=0.03$ is about an order of magnitude smaller than the value from our \nmodel. This discrepancy cannot be resolved, even by varying the turnover time\nof the largest eddies $T_{e}$ (the dashed line). One possible cause for the\ninconsistency is particle clumping and concentration in turbulent eddies,\nwhich are not included in our model. The simulation shows strong\nparticle clumping when $Z=0.03$ in the R10 run and also temporal clumping for\n$Z=0.03$ in the R30 run. Such clumping of particles in turbulent eddies may\nsuppress diffusion of particles compared to the no-clumping cases of\n$Z=0.01$ and could be a cause of a decrease in $\\alpha$ when\n$Z$ is increased in the simulation. Even if the\nlargest particles (of $T_{s}\\sim1$) concentrate in clumps, the smallest\nparticles (of $T_{s}\\la 0.1$) do not clump, and continue to follow\nthe turbulent diffusion of the gas (Fig. 7 of Bai \\& Stone 2010). In Figure\n\\ref{fig:alpha-bai}$b$, the turbulent diffusion coefficient in the ``radial\ndirection'' of the smallest particles in the simulation, $D_{x}$, is compared\nwith the ``vertical'' diffusion coefficient $\\alpha$ in our model. Note that\nwe compare diffusion coefficients in the different\ndirections. Since the smallest particles spread out to high altitudes where\nturbulence is weak, it is difficult to measure the vertical diffusion\ncoefficient for the smallest particles in the simulation. In Figure\n\\ref{fig:alpha-bai}$b$, though we still see a discrepancy compared with\nthe R30 run, our model appears more consistent with the simulation\nresults, suggesting that our model properly predicts the ``gas''\ndiffusion coefficient. \n\n\\section{DISCUSSION}\n\\label{sec:discuss}\n\n\\subsection{The Radial Drift Velocity and Collision Velocity of Dust\nParticles}\n\\label{sec:dustvelocity}\n\n\\begin{figure}[ptb]\n\\epsscale{1.1} \\plotone{f10.eps} \\caption{The radial drift velocity, $\\bar\n{v}_{d,r}$ (solid line), and the relative velocity of particles due to\nturbulence, $\\Delta v_{12}$ (dashed line). The contributions of individual\ndrag, $\\bar{v}_{\\mathrm{ind},r}$, and collective drag, $\\bar{v}_{\\mathrm{col\n,r}$, to the radial drift velocity are plotted with dotted lines.\n\\label{fig:velocity\n\\end{figure}\n\nIf the dust-to-gas ratio in the dust layer were larger than unity, the radial\ndrift velocity of dust particles would be lower than the value that the\nparticles would have in a gas-rich environment because the gas drag\nforce could not accelerate sufficiently against the large inertia of the\ndust. This effect was pointed out by NSH86 and is seen in equation\n(\\ref{eq:vdr1}) for the radial velocity due to individual drag. If the\ndust-to-gas ratio were much larger than unity, dust accretion would\nbe caused by collective drag exerted from the slower-orbiting\nupper gas layer (Weidenschilling 2003). The radial drift velocity of the\ndust is thus a function of the turbulence strength, $\\alpha$. For\nweaker turbulence, the \ndust-to-gas ratio in the dust layer is higher, and the individual drag is\nweaker. The collective drag is also weaker at smaller $\\alpha$ because the\nReynolds stress $P_{\\theta z}$ is proportional to $\\alpha$. The relative\nvelocity (or collision velocity) of the dust particles is also a function of\n$\\alpha$. The radial drift and collision velocities are the important factors in\nthe dust growth process. In the previous sections, the turbulence strength\n$\\alpha$ and the dust layer thickness $h_{d}$ have been determined\nself-consistently. Using these results, the radial drift velocity and collision\nvelocity of the dust particles are estimated.\n\nThe radial drift velocity is calculated separately for the components due to\nindividual drag and due to collective drag. For each component, the radial\ndrift velocity is averaged in the vertical direction. First, the averaged\nvalue of the radial drift velocity due to individual drag is calculated\nfrom equation (\\ref{eq:vdr1}),\n\\begin{eqnarray}\n\\overline{v}_{\\mathrm{ind},r} &=& \\frac{1}{\\Sigma_{d}}\\int_{-\\infty}^{\\infty\n\\rho_{d}v_{d,r}dz \\nonumber \\\\\n&=& -2\\eta v_{\\mathrm{K}}T_{s}\\frac{1}{\\Sigma_{d}}\\int_{-\\infty}^{\\infty\n}\\rho_{d}\\frac{1}{T_{s}^{2}+\\beta^{2}}dz~.\n\\end{eqnarray}\nThe radial drift velocity due to collective drag is calculated from the\nvertically-averaged angular momentum loss of the dust component. Integrating\nequation (\\ref{eq:dlddt}) gives\n\\begin{equation}\n\\frac{dL_{d}}{dt}=\\int_{-\\infty}^{\\infty}\\frac{\\partial l_{d}}{\\partial\nt}dz=-2\\eta v_{\\mathrm{K}}^{2}T_{s}\\Sigma_{d,\\mathrm{vis}}~,\n\\end{equation}\nwhere $\\Sigma_{d,\\mathrm{vis}}$ is given by equation (\\ref{eq:sig_vis2}). This\nangular momentum loss causes a radial drift velocity $\\bar{v}_{\\mathrm{col},r}$ given by\n\\begin{equation}\n\\bar{v}_{\\mathrm{col},r}=\\frac{2}{v_{\\mathrm{K}}\\Sigma_{d}}\\frac{dL_{d}\n{dt}=-4\\eta v_{\\mathrm{K}}T_{s}\\frac{\\Sigma_{d,\\mathrm{vis}}}{\\Sigma_{d}}~.\n\\end{equation}\nThe total radial drift velocity is $\\bar{v}_{d,r}=\\overline{v}_{\\mathrm{ind\n,r}+\\bar{v}_{\\mathrm{col},r}$, and is shown in Figure \\ref{fig:velocity}\nas a solid line, for the case in which the particle size is chosen to maximize the\nradial velocity ($T_{s}=1$). In plotting this figure, we adopt\nthe model parameters at 1 AU of the minimum-mass-solar-nebula model of\nHayashi (1981): $h_{g}\/r=0.0326$, $\\eta=1.80\\times10^{-3}$, and\n$\\tilde{\\eta}=2.92 \\times 10^{-3}$. For a small disk metallicity\n$Z_{\\mathrm{disk}}$, the radial drift velocity is as large \nas 50 m s$^{-1}$, and it decreases with $Z_{\\mathrm{disk}}$. For $Z_{\\mathrm{disk}}>0.08$, the radial velocity is dominated by collective drag (see the dotted\nlines) as pointed out by Weidenschilling (2003). For such large\n$Z_{\\mathrm{disk}}$, the radial drift velocity due to collective drag also\ndecreases with $Z_{\\mathrm{disk}}$, and then it becomes as small as $1$ m s$^{-1}$\nfor $Z_{\\mathrm{disk}}=0.2$. The radial drift velocity is strongly suppressed\nfor a sufficiently massive dust layer.\n\nThe relative velocity of dust particles due to turbulence, $\\Delta v_{12}$, is\ncalculated for $T_{s}=1$ by substituting $\\mathrm{St}_{1}=1$ and\n$\\mathrm{St}_{2}=0$ into equation (29) in Ormel \\& Cuzzi (2007), and is shown\nin Figure \\ref{fig:velocity} as a dashed line. The collision velocity is\nestimated by the larger of $\\Delta v_{12}$ and $\\bar{v}_{d,r}$. For small\ndisk metallicities $Z_{\\mathrm{disk}}$, the collision velocity is dominated by\nthe radial drift and is as large as 50 m s$^{-1}$, while for large $Z_{\\mathrm{disk}}>0.02$, it is dominated by turbulence. The maximum value of the collision\nvelocity due to turbulence is about 30 m s$^{-1}$, and it decreases with\n$Z_{\\mathrm{disk}}$ for $Z_{\\mathrm{disk}}>0.03$. Thus, if the dust\nparticles could survive collisions of $30$ m s$^{-1}$, as suggested by the numerical\nsimulation of collisions of dust aggregates (Wada et al. 2010) and\n$Z_{\\mathrm{disk}}\\ga0.03$, the dust particles would be able to grow without\nbeing reduced to small fragments.\n\n\\subsection{Radial Dependence of the Midplane Dust-to-Gas Ratio}\n\\label{sec:disketa}\n\n\\begin{figure}[ptb]\n\\epsscale{1.1} \\plotone{f11.eps} \\caption{The midplane dust-to-gas\n ratio $f_{\\mathrm{mid}}$ is plotted against $\\tilde{\\eta}$ for various values\n of the disk metallicity $Z_{\\rm disk}$. The top axis shows the\n corresponding radii in the gas disk model by Hayashi (1981), in which\n the disk temperature is $T = 278 (r\/{\\rm AU})^{-1\/2}$K. The dust\n particles are assumed to be small ($T_s \\ll 1$) such that\n $f_{\\mathrm{mid}}$ is independent of $T_s$.}\n\\label{fig:fmid-eta\n\\end{figure}\n\nAs shown in \\S\\ref{sec:turbstr}, the turbulence strength and the dust\nlayer structure depend on properties of the gas disk only through\n$\\tilde{\\eta}$. For larger $\\tilde{\\eta}$, the accretion velocity of the\ndust is faster, and consequently the dust layer is thicker due to\nstronger turbulence. Figure \\ref{fig:fmid-eta} shows how the midplane\ndust-to-gas ratio varies with $\\tilde{\\eta}$. The figure is plotted in the limit\nof $T_s \\ll 1$. For such small particles, $f_{\\rm mid}$ is\nindependent of $T_s$ (see Fig. \\ref{fig:fdmid}). The midplane\ndust-to-gas ratio $f_{\\rm mid}$ decreases with $\\tilde{\\eta}$ as shown\nin equation (\\ref{eq:fdmid1}). For a gas disk model with a power-law temperature\nprofile, $T \\propto r^{-q}$, $\\tilde{\\eta}$ behaves as\n$\\tilde{\\eta} \\propto r^{1-q}$. The top axis of Figure\n\\ref{fig:fmid-eta} indicates corresponding locations in the disk\nmodel by Hayashi (1981), i.e., $T = 278 (r\/{\\rm AU})^{-1\/2}$K. In a\ndisk with the standard value of the disk metallicity $Z_{\\rm disk}=0.01$,\n$f_{\\rm mid}$ is less than unity in the most part of the disk except $r\n\\la 0.1$AU. In disks with $Z_{\\rm disk}=0.3$, $f_{\\rm mid}$\nexceeds $100$ for $r \\la 1$AU. We discuss the condition for planetesimal\nformation through gravitational instability of the dust layer. In a gas\ndisk with a surface density profile $\\Sigma_g \\propto r^{-p}$, the\nmidplane gas density scales as $\\rho_g \\propto r^{-p+q\/2-3\/2}$. The\ncondition that the dust density exceeds the Roche density ($\\rho_{\\rm R}\n\\propto r^{-3}$) is $f_{\\rm mid} > \\rho_{\\rm R}\/\\rho_g \\propto\nr^{p-q\/2-3\/2} \\propto r^{p-7\/4}$ in a disk model with\n$q=1\/2$. From the lower line of equation (\\ref{eq:fdmid1}), $f_{\\rm\n mid}$ decreases as $f_{\\rm mid} \\propto r^{-(1-q)\/(2-\\delta)}\\propto\nr^{-0.47}$. If $p>1.3$, the inner part of the disk is more suitable\nfor gravitational instability, and vice versa.\n\n\\subsection{Does the Liberated Gravitational Energy Go into\nTurbulence?}\n\\label{sec:ene-into-tur}\n\nIn this paper, we calculate the liberated gravitational energy from dust\naccretion, assuming some fraction of the liberated energy is transferred to\nturbulence. The estimate of the dust accretion rate is based on the formula for the\nparticle terminal velocity derived in NSH86. ``The terminal velocity'' means\nthat all the liberated energy is consumed by gas drag, converting directly into\nthe thermal energy of the dust particles and of the surrounding gas molecules.\nThus, one may expect that only a small fraction (or nothing) of the liberated\nenergy would be used for maintaining turbulence. However, a comparison with\nthe simulation of KH instability by JHK06 shows that the efficiency factor\n$C_{\\mathrm{eff}}\\approx0.2$ is not negligibly small. In the following subsection, we\ndiscuss the validity of using the particle terminal velocity for calculating\nthe energy supply rate to turbulence. The energy liberation rate from the\naccreting dust calculated in \\S \\ref{sec:eneinp} is compared with the deposit\nrate of the free energy that is the source of several instabilities, such as KH\ninstability and streaming instability.\n\n\\subsubsection{Kelvin-Helmholtz Instability}\n\nThe free energy that induces the KH instability originates from the velocity\ndifference between the midplane dust layer and the upper gas layer, and is\nstored as the dust particles settle to the midplane. We estimate the deposit\nrate of the free energy during dust sedimentation, and show that it has the\nsame order of magnitude as the energy liberation rate from the dust accretion\ntowards the star. Consider two states of dust distribution: the\ninitial state, in which the dust particles are distributed uniformly in the gas\ndisk, and the final state, in which all the dust has settled at\nthe midplane. In the initial state, there is no vertical shear in the disk, and in the final state, the velocity difference $\\Delta\nv_{\\theta}=\\eta v_{\\mathrm{K}}$ appears between the midplane dust layer and\nthe upper gas layer. The free energy for KH instability is $\\Delta\nE_{\\mathrm{KH}}\\sim\\frac{1}{2}\\Sigma_{d}\\Delta v^{2}\\sim\\frac{1}{2}\\Sigma\n_{d}\\eta^{2}v_{\\mathrm{K}}{}^{2}$, for $\\Sigma_{d}\\ll\\Sigma_{g}$. The settling\ntimescale is $\\tau_{\\mathrm{sed}}=(T_{s}^{2}+1)\/(T_{s} \\Omega_{\\mathrm{K}})$,\nand then the energy deposit rate is\n\\begin{equation}\n\\frac{\\Delta E_{\\mathrm{KH}}}{\\tau_{\\mathrm{sed}}}\\sim\\frac{1}{2}\\eta\n^{2}v_{\\mathrm{K}}^{2} \\Omega_{\\mathrm{K}} \\frac{T_{s}}{T_{s}^{2}+1} \\Sigma_{d} \\ ,\n\\end{equation}\nwhich is the same order as the energy liberation rate of the accreting dust\n(eqs. [\\ref{eq:ene_drag}] and [\\ref{eq:sig_drag2}]). Hence, the deposition rate\nof the free energy for KH instability can be estimated by the energy\nliberation rate of the accreting dust.\n\n\\subsubsection{Streaming Instability}\n\\label{sec:streaminst}\n\nThe free energy for streaming instability originates from the velocity\ndifference between the dust particles and the surrounding gas. When streaming\ninstability begins, the velocity difference decreases as the free energy is\nconsumed by inducing turbulence. In fact, this decrease in the velocity difference\ncan be seen even in the linear growth regime. Youdin \\& Goodman (2005)\nshowed in their Figure 6 that the velocity difference between the dust and gas\ndecreases (increases) at the locations where the particle density increases\n(decreases). The spatially averaged value of the free energy decreases as the\nperturbation grows. Thus, without an energy supply, streaming instability would\ncease. Given a state in which the velocity difference between the dust and the\ngas has reduced, the dust particles are no longer in equilibrium: the\ngravity, the centrifugal force, and the gas drag force are not in balance. The\ndust particles are accelerated and the velocity difference from the gas rises\nagain. To estimate the effect of the energy deposition on the velocity\ndifference, we consider a state in which the dust density is similar to the\ngas density, $\\rho_{d}\\sim\\rho_{g}$. In such a state, streaming\ninstability occurs efficiently with a growth time of the order of the\nKeplerian time (for the short wave branch, Youdin \\& Goodman 2005; Youdin \\&\nJohansen 2007). For dust particles of $T_{s}<1$, the terminal velocities of\nthe dust and of the gas are of the order of $T_{s}\\eta v_{\\mathrm{K}}$ (eq.\n[\\ref{eq:vdr}], [\\ref{eq:vgr}]), and thus, the free energy per unit area is\n$\\Delta E_{\\mathrm{str}}\\sim\\Sigma_{d}T_{s}^{2}\\eta^{2}v_{\\mathrm{K}}^{2}$. This\ndeposition of free energy occurs during the acceleration phase of the dust,\nand thus in the stopping time $T_{s}\\Omega_{\\mathrm{K}}^{-1}$, and then the energy is\ntransferred to turbulence in the growth timescale $\\Omega_{\\mathrm{K}}^{-1}$ of\nstreaming instability. Thus, the timescale for the energy deposit in\nturbulence is the sum of these timescales, and for $T_{s}\\la1$, $\\tau\n_{\\mathrm{str}}=T_{s}\\Omega_{\\mathrm{K}}^{-1}+\\Omega_{\\mathrm{K}}^{-1}\\sim\\Omega_{\\mathrm{K}}^{-1}$. The\nenergy deposition rate is\n\\begin{equation}\n\\frac{\\Delta E_{\\mathrm{str}}}{\\tau_{\\mathrm{str}}}\\sim\\eta^{2}\nv_{\\mathrm{K}}^{2} \\Omega_{\\mathrm{K}} T_{s}^{2}\\Sigma_{d} \\ ,\n\\end{equation}\nwhich is smaller by a factor $T_{s}$ than the estimate from the dust accretion\nrate (eq.[\\ref{eq:ene_drag}]). Hence, our estimate, derived from the dust accretion\nrate, is appropriate for particles of $T_{s}\\sim1$. For smaller particles\n($T_{s} \\ll1$), the energy deposition rate is higher for KH instability than for\nstreaming instability, and KH instability is expected to operate first. The\nenergy deposition rate for instability (either for KH or streaming instabilities)\nis estimated from the energy liberation rate due to dust accretion.\n\n\\section{SUMMARY}\n\nIn this paper, we discuss turbulence induced in the dust layer. The\n turbulence strength or the parameter $\\alpha$ is determined using the energetics of\ndust accretion towards the central star. The key concept is that the dust\nparticles reside in a deeper potential than the gas. The effective potential,\nincluding the gas pressure, is $-GM(1-2\\eta)\/r$ for the gas,\nand $-GM\/r$ for the dust. When angular momentum is transferred\nfrom the dust to the gas through gas drag, the dust particles lose more energy\nthan the gas gains. The excess energy can be used for exciting\nturbulence. If the dust accretion due to gas drag is a primary source of energy\nliberation, i.e., if the gas accretion rate due to turbulence is much smaller\nthan the dust accretion rate, then the turbulence strength is determined by the\nenergy supply rate from the dust accretion. This is not the case if the gas\ndisk itself is turbulent via, e.g., MRI. If the dust layer is composed of large\nparticles with stopping time $T_{s}\\gg1$, then the gas accretion may\ndominate the dust accretion, as discussed in \\S \\ref{sec:coldrag}.\n\n We estimate the dust accretion rate using the terminal velocity profiles of the\ndust particles in a laminar disk derived by NSH86. The expected turbulence\nstrength and corresponding structure of the dust layer from our analysis agree\nwith the previous analytical result on the marginally KH-unstable dust layer\nby S98. As our analysis does not assume tight coupling of the dust to the\ngas, nor specify the mechanism of instability that induces turbulence, it is\nconsidered an extension of the analysis of S98 to a more general physical\nsituation of the dust layer. The results of this paper agree with\nthe results in Michikoshi \\& Inutsuka (2006), which analyzes KH instability\nof the dust layer composed of particles with large stopping times ($T_{s}>1$),\nas shown in Figure \\ref{fig:michi}. \n\nOur analysis shows that, for particles of\n$T_{s}\\la1$, the turbulence strength is smaller than $\\alpha_{\\max}\\sim\nC_{\\mathrm{eff}}{\\tilde{\\eta}}T_{s}$, where $C_{\\mathrm{eff}}=0.19$ is the\nefficiency of the energy supply to turbulence (see Fig. \\ref{fig:alpha} and\neq. [\\ref{eq:alpha2}]). This strength reaches a maximum when the disk\nmetallicity is $Z_{\\mathrm{disk}}\\sim\\sqrt{C_{\\mathrm{eff}}\\tilde{\\eta\n}}\\sim10^{-2}$. Modifying the disk metallicity from the standard value,\n$10^{-2}$, by any process, results in weaker turbulence and a thinner dust\nlayer, and consequently may accelerate the growth process of the dust\nparticles, as pointed out in S98.\n\nComparison of our results with previous numerical simulations of KH and streaming\ninstabilities by JHK06 and Bai \\& Stone (2010) shows quantitative agreement\nwith our analysis for dust particles of $T_{s}\\la0.1$, although there may be a\nqualitative disagreement for $T_{s}\\ga1$ particles (see Figs.\n\\ref{fig:alpha-joh} $-$ \\ref{fig:alpha-bai}). Hence, we conclude that\nturbulence in the dust layer is controlled by the energy supply from the dust\naccretion due to gas drag, provided that the dust particles are not so large\nthat $T_{s}\\ga1$. In such a layer, turbulence strength is estimated by the dust\naccretion rate (eq.[\\ref{eq:alpha}]).\n\n\\acknowledgements\nThis work was stimulated by discussions with Minoru Sekiya. We are\ngrateful to Anders Johansen for providing detailed information on his\nsimulations, and to Chris Ormel and Eugene Chiang for useful\ndiscussions. We also thank an anonymous referee for helpful comments.\nThis work was supported in part by Grants-in-Aid for Scientific Research,\nNos. 20540232, 22$\\cdot$2942, and 22$\\cdot$7006 from the Ministry of\nEducation, Culture, Sports, Science, and Technology (MEXT), Japan.\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Statement of the main results}\n\\label{sec:main-statement}\n\n\\subsection{Setup}\n\nLet $I:=[0,1] \\subset {\\mathbb R}$ and for $k, n \\in {\\mathbb N}} \\def\\Q{{\\mathbb Q}} \\def\\P{{\\mathbb P},\\ k \\ge n$,\ndenote by $\\pi^k_n:I^k\\to I^n$ the projection map. We sometimes omit\n$k$ if its meaning is clear from the context.\n\nPfaffian functions, introduced by Khovanskii in \\cite{khov:dan,\n khov:book}, are analytic functions satisfying triangular systems of\nPfaffian (first order partial differential) equations with polynomial\ncoefficients. We refer the reader to\n\\cite{gv:complexity-computations} for precise definition of Pfaffian\nfunctions, based on {\\em Pfaffian chain}, and examples of Pfaffian\nfunctions in an open domain $G \\subset {\\mathbb R}^k$, which we assume\nhere for simplicity to be given by a product of intervals.\n\n\\begin{Def}[semi-Pfaffian set]\n Let $G $ be an open set in ${\\mathbb R}^k$ and $I^k \\subset G$. A\n set $X \\subset I^k$ is called \\emph{(restricted) semi-Pfaffian} if\n it consists of points in $I^k$ satisfying a Boolean combination of\n atomic equations and inequalities of the kind $f=0$ or $f>0$, where\n $f$ is a Pfaffian functions defined in $G$. The \\emph{format} of $X$\n is the number of variables $k$ and the {\\em degree} of $X$ is the\n sum of degrees of all the Pfaffian functions appearing in the atomic\n formulas (i.e. the degrees of the polynomials defining these\n Pfaffian functions).\n\\end{Def}\n\nNote that the degree of $X$ bounds from above the number of all atomic\nequations and inequalities. We assume that a Pfaffian chain has been\nfixed once and for all, and all semi-Pfaffian sets under consideration\nare defined from this single Pfaffian chain.\n\n\\begin{Def}[sub-Pfaffian set]\\label{def:sub-pfaff}\n A set $Y \\subset I^n$ is called {\\em (restricted)\n sub-Pfaffian}\\footnote{Below we will consider only\n \\emph{restricted} sub-Pfaffian sets, and refer to them simply as\n sub-Pfaffian.} if $Y = \\pi^k_n (X)$ for a semi-Pfaffian set\n $X \\subset I^k$.\n\\end{Def}\n\nIn the special case of a {\\em semi-algebraic set} $X$, the\nTarski-Seidenberg theorem states that the set $Y = \\pi^k_n (X)$ is\nalso semi-algebraic, i.e., is a set of points satisfying a Boolean\ncombination of polynomial equations and inequalities. By contrast, a\nsub-Pfaffian set may not be semi-Pfaffian. \n\nIt is customary in the literature to define the format and degree of a\nsub-Pfaffian set as in Definition~\\ref{def:sub-pfaff} to be the format\nand degree of the set $X$. Below we introduce a variant of these\nnotions which turns out to behave better with respect to cell\ndecompositions. To avoid confusion we refer to these modified notions\nas *-format and *-degree. See Remark~\\ref{rem:comparing-formats} for a\ncomparison between the starred and unstarred versions of format and\ndegree.\n\n\n\\begin{Def}[*-format and *-degree of projections]\\label{def:format}\n If $\\{X_\\alpha\\},\\ X_\\alpha \\subset I^{k_\\alpha}$ is a finite\n collection of semi-Pfaffian sets and $X_\\alpha^\\circ$ is a connected\n component of $X_\\alpha$, we define the *-format of the sub-Pfaffian\n set\n $Y:=\\bigcup_\\alpha \\pi^{k_\\alpha}_n(X_\\alpha^\\circ) \\subset I^n$,\n with this particular representation as a union of projections, to be\n the maximum among the formats of sets $X_\\alpha$, and the *-degree\n of $Y$ to be the sum of the degrees of sets $X_\\alpha$.\n\\end{Def}\n\nStrictly speaking, the *-format $\\cF$ and *-degree $D$ is associated\nto a particular presentation of a sub-Pfaffian set in the form\nprescribed by Definition~\\ref{def:format}, with different\npresentations possibly giving rise to different pairs $(\\cF,D)$. By a\nslight abuse of notation we will say that a sub-Pfaffian set has\nformat $\\cF$ and degree $D$ if there exists \\emph{some} presentation\nhaving this data.\n\nWe remark that since the restricted sub-Pfaffian sets form an\no-minimal structure, the connected components of semi-Pfaffian (or\neven sub-Pfaffian) sets, as well as their projections, are again\nsub-Pfaffian, so the sets $Y$ for which *-format and *-degree are\nintroduced in Definition~\\ref{def:format} are indeed sub-Pfaffian.\n\n\\begin{Rem}\n It may be useful below for the reader to consider the notions of\n *-format and *-degree as sub-Pfaffian analogs of the notions of\n dimension and degree in the theory of algebraic or semi-algebraic\n geometry. Our main objective is to obtain, as in the semialgebraic\n case, bounds that depend polynomially on the degree for a fixed\n format.\n\\end{Rem}\n\nIn what follows we write that $a$ is $\\const(b)$ (resp., $a$ is\n$\\operatorname{poly}} \\def\\J{\\operatorname{J}_b(c)$) as shorthand notation meaning that $a \\le \\gamma(b)$\n(resp., $a \\le (c+1)^{\\gamma(b)}$) where $\\gamma(\\cdot)$ is some\n\\emph{universally fixed} function (which may be different for each\noccurrence of this notation in the text). Here $a$ and $c$ denote\nnatural numbers, and $b$ can involve one or several arguments of any\ntype. We also allow several arguments in $\\operatorname{poly}} \\def\\J{\\operatorname{J}_b(c_1,\\ldots,c_n)$\nwhich we interpret as $\\operatorname{poly}} \\def\\J{\\operatorname{J}_b(c_1+\\cdots+c_n)$. For each occurrence\nof an asymptotic notation in this text, the implied functions can be\neffectively and explicitly computed from the data defining the\nPfaffian chain (degrees of the differential equations involved), and\nwe omit these computations in the interest of brevity.\n\n\n\\subsection{Cell decompositions}\n\nWe recall the standard definitions of a cylindrical cell and a\ncylindrical cell decomposition. Later in the paper we consider only\ncylindrical decompositions and omit the prefix ``cylindrical'' for\nbrevity.\n\n\\begin{Def}[Cylindrical cell] A \\emph{cylindrical cell} is defined by\n induction as follows.\n \\begin{enumerate}\n \\item Cylindrical 0-cell in ${\\mathbb R}^n$ is an isolated point.\n \\item Cylindrical 1-cell in ${\\mathbb R}$ is an open interval\n $(a,b) \\subset {\\mathbb R}$.\n \\item For $n \\ge 2$ and $0 \\le \\ell < n$ a cylindrical\n $(\\ell+1)$-cell in ${\\mathbb R}^n$ is either a graph of a\n continuous bounded function $f:\\ C \\to {\\mathbb R}$, where $C$ is\n a cylindrical a cylindrical $(\\ell+1)$-cell in\n ${\\mathbb R}^{n-1}$, or else a set of the form\n \\begin{multline}\n \\{(x_1, \\ldots ,x_n) \\in {\\mathbb R}^n : \\ (x_1, \\ldots , x_{n-1}) \\in C \\\\\n \\text{ and } f(x_1,\\ldots,x_{n-1})s_\\beta\n \\end{align}\n holds uniformly over $C$. Indeed, otherwise there should be a point\n $c_0\\in C$ where $s_\\alpha(c_0)=s_\\beta(c_0)$ and a neighborhood\n where this does not hold identically. This implies that\n $(c_0,c_\\alpha(c_0))$ belongs to a stratum of dimension at most\n $(k-1)$, or to the frontier of a $k$-dimensional stratum of\n $Z_{\\alpha,\\beta}$, and hence $c_0\\not\\in G$ contradicting\n $C\\subset G$.\n\n By the above, the set of sections obtained from any of the\n $X_\\alpha\\in\\Pi_k$ is\n $$\\{s_1<\\cdots d_{FD}(f,g)$.\nBy definition of the functional distortion metric, there are maps\n \\begin{equation*}\n \\begin{tikzcd}\n \\X \\ar[start anchor=north east, end anchor=north west]{r}{\\phi}\n & \\Y \\ar[start anchor=south west, end anchor=south east]{l}{\\psi}\n \\end{tikzcd}\n \\end{equation*}\nwhich satisfy the requirements of Definition~\\ref{Def:FuncDistortion}, in particular, $\\max\\{\\|f-g\\phi\\|_\\infty, \\|g-f\\psi\\|_\\infty\\} \\leq \\e$.\nThen it is obvious that $\\phi(f\\inv(I)) \\subset g\\inv(I^\\e)$ and $\\psi(g\\inv(I)) \\subset f\\inv(I_\\e)$ for all intervals $I \\subset \\R$.\nThis implies that there is a map induced by $\\phi$\n\\begin{equation*}\n \\begin{array}{rccc}\n \\phi^*_I: & \\pi_0(f\\inv(I)) & \\longrightarrow & \\pi_0(g\\inv(I^\\e))\n \\end{array}\n\\end{equation*}\nand a map induced by $\\psi$\n\\begin{equation*}\n \\begin{array}{rccc}\n \\psi^*_I: & \\pi_0(g\\inv(I)) & \\longrightarrow & \\pi_0(f\\inv(I^\\e))\n \\end{array}\n\\end{equation*}\nfor every $I$.\nBy the functoriality of $\\pi_0$, these maps commute with the maps induced by the inclusions $\\iota: I \\hookrightarrow J$ and $\\iota_\\e: I^\\e \\hookrightarrow J^\\e$; that is, the diagrams\n\\begin{equation*}\n \\begin{tikzcd}\n \\pi_0(f\\inv(I)) \\ar{r}{\\phi^*_I} \\ar{d}{\\iota^*}& \\pi_0(g\\inv(I^\\e)) \\ar{d}{\\iota_\\e^*} &&\n \\pi_0(g\\inv(I)) \\ar{r}{\\psi^*_I} \\ar{d}{\\iota^*}& \\pi_0(f\\inv(I^\\e)) \\ar{d}{\\iota_\\e^*}\\\\\n \\pi_0(f\\inv(J)) \\ar{r}{\\phi^*_J}& \\pi_0(g\\inv(J^\\e)) &&\n \\pi_0(g\\inv(J)) \\ar{r}{\\psi^*_J}& \\pi_0(f\\inv(J^\\e)) \\\\\n \\end{tikzcd}\n\\end{equation*}\ncommute.\nIn addition, the functoriality of $\\pi_0$ implies that the diagrams\n\\begin{equation*}\n \\begin{tikzcd}\n \\pi_0f\\inv(I) \\ar{r}\\ar{dr} & \\pi_0 g\\inv(I^\\e) \\ar{d} && \\pi_0g\\inv(I) \\ar{r} \\ar{dr} & \\pi_0 f\\inv(I^\\e) \\ar{d} \\\\\n & \\pi_0f\\inv(I^{2\\e}) && & \\pi_0g\\inv(I^{2\\e}) \n\\end{tikzcd}\n\\end{equation*}\ncommute.\nThese are exactly the properties necessary to call $\\phi^*$ and $\\psi^*$ an $\\e$-interleaving of the associated cosheaves. Since the above holds for any $\\e > d_{FD}(f,g)$, we conclude $d_I(f,g) \\leq \\e = d_{FD}(f,g)$.\n\\end{proof}\n\n\n\\subsection{The Hard Direction}\n\nIn order to show $d_{FD}((\\X,f), (\\Y,g)) \\leq d_I((\\X,f), (\\Y,g)) $, we need to start with an $\\e$-interleaving, $\\phi:(\\X,f) \\to \\TT_\\e(\\Y,g)$ and $\\psi:(\\Y,g) \\to \\TT_\\e(\\X,f)$, and construct a pair of maps satisfying the requirements of the functional distortion distance.\nTo do this, note that the map $\\phi$ induces a multimap $\\bar \\phi: \\X \\to \\Y_\\e$, which sends a point $x$ to the entire equivalence class of $\\phi(x)$, thought of as a subset of $\\Y_\\e$.\nConcretely, letting $q:\\Y_\\e \\to \\TT_\\e(\\Y)$ denote the Reeb graph quotient map, we have $\\bar\\phi = q^{-1}\\circ \\phi$.\nFor any $\\delta$, we can construct the multimap, $\\bar\\phi_\\delta:\\X \\to \\Y_\\e$ which sends $x$ to $\\bar\\phi(B_\\delta(x))$, where $B_\\delta(x) = \\{x' \\mid d_f(x,x')<\\delta\\}$. Explicitly, we have\n\\begin{equation*}\n\\bar\\phi_\\delta(x) = \\{(y',t') \\in \\Y_\\e \\mid (y',t') \\in \\bar\\phi(x'), d_f(x,x')<\\delta\\},\n\\end{equation*}\nSee Fig.~\\ref{Fig:DeltaMap} for an example.\n\n\\begin{figure}\n \\centering\n \\includegraphics[scale = .8]{SmoothingReebGraph-DeltaMap.pdf}\n \\caption{An example for determining the map $\\bar\\phi_\\delta$. Given the red point $x \\in \\X$, the red solid region in $\\X$ is $B_\\delta(x)$. Then we can look at $\\phi(B_\\delta(x))$, the red region in $\\TT_\\e(\\Y)$. The set $\\bar\\phi_\\delta(x)$ in $\\Y_\\e$ consists of the points which map into $\\phi(B_\\delta(x))$ in $\\TT_\\e(\\Y)$ under the quotient map $q$. }\n \\label{Fig:DeltaMap}\n\\end{figure}\n\nWe want to show that the map $\\bar \\phi_\\delta: \\X \\to \\Y_\\e$ satisfies the assumptions of Theorem~\\ref{Thm:Selection}. \n\\begin{enumerate}\n \\item \n Since $\\X$ is a finite graph, it is compact and thus trivially paracompact.\nIn addition, because it is a graph, it has covering dimension 1. \n\\item\nWe can embed the graph $\\Y$ in $\\R^3$. \nIn this context, it is a closed subset of a completely metrizable space, and thus is completely metrizable.\nTherefore, $\\Y_\\e$ is also completely metrizable since it is the product of two completely metrizable spaces.\n\\item To show that $\\bar\\phi_\\delta$ is LSC, let $U \\subset \\Y_\\e$ open and $x \\in \\bar\\phi_\\delta\\inv(U)$ given.\n This means that there is an $x' \\in B_\\delta(x)$ such that $\\bar\\phi(x')\\cap U \\neq \\emptyset$.\n Let $r>0$ be a radius such that $B_r(x')\\subset B_\\delta(x)$. \n In particular, this means that $d_f(x,x')<\\delta-r$.\n We now want to show that $B_r(x) \\subset \\bar\\phi_\\delta\\inv(U)$.\n Let $x'' \\in B_r(x)$. \n We know that $x' \\in B_\\delta(x'')$ since \n \\begin{equation*}\n d_f(x',x'') \\leq d_f(x',x) + d_f(x,x'') d_I(f,g)$ and for any $\\delta>0$, this completes the proof.\n\\end{proof}\n\nPutting together Lemmas \\ref{Lem:Inter leq FD} and \\ref{Lem:FD leq Inter}, our main theorem is immediate.\n\\begin{theorem}\n\\label{Thm:Main}\n The functional distortion metric and the interleaving metric are strongly equivalent. That is, given any Reeb graphs $(\\X,f)$ and $(\\Y,g)$,\n \\begin{equation*}\n d_I(f,g) \\leq d_{FD}(f,g) \\leq 7d_I(f,g).\n \\end{equation*}\n\\end{theorem}\n\n\n\n\n\\section{Relationship Between the Interleaving and Bottleneck Distances}\n\\label{sec:bottleneck}\n\n\nHaving strongly equivalent metrics means that we can quickly pass back and forth many of the properties associated to the metrics.\nFor example, the bottleneck stability bound for persistence diagrams in terms of the functional distortion distance~\\cite{Bauer2014} says the following (for the definitions of the persistence diagrams $\\Dg_0(f), \\ExDg_1(f)$ associated to a function~$f$ and of the bottleneck distance $d_B$ we refer the reader to~\\cite{EH09}):\n\\begin{theorem}\nGiven two Reeb graphs $(\\X,f)$ and $(\\Y,g)$,\n\\begin{equation*}\n d_B(\\Dg_0(f),\\Dg_0(g)) \\leq d_{FD}(f,g)\n\\end{equation*}\n and \n\\begin{equation*}\n d_B(\\ExDg_1(f),\\ExDg_1(g)) \\leq 3d_{FD}(f,g).\n\\end{equation*}\n\\end{theorem} \n \nCombining this result with Theorem~\\ref{Thm:Main} gives an immediate stability result relating the interleaving distance with the bottleneck distance.\n\\begin{theorem}\nGiven two Reeb graphs $(\\X,f)$ and $(\\Y,g)$,\n\\begin{equation*}\n d_B(\\Dg_0(f),\\Dg_0(g)) \\leq 7d_I(f,g)\n\\end{equation*}\n and \n\\begin{equation*}\n d_B(\\ExDg_1(f),\\ExDg_1(g)) \\leq 21d_I(f,g).\n\\end{equation*}\n \n\\end{theorem}\n\n\n\n\n\n\\section{Discussion}\n\nIn this paper, we study the relation between the two existing distance measures for Reeb graphs, and show that they are strongly equivalent on the space of Reeb graphs. \nThis relationship will be a powerful tool for understanding convergence properties of the different metrics. \nFor example, if we have a Cauchy sequence in one metric, we have a Cauchy sequence in the other and can therefore pass around completeness results. \nThis relationship also means that algorithms for computation and approximation of the metrics can be written using whichever method is most helpful and applicable to the context.\n\n\nThese two distances in general may not be the same. An immediate question is whether the relations provided in Theorem \\ref{Thm:Main} are tight. \nIn particular, it is easy to construct examples where the bound $d_I(f,g) \\leq d_{FD}(f,g)$ of Lemma~\\ref{Lem:Inter leq FD} is tight; it will be interesting to investigate the tightness of the bound $d_{FD}(f,g) \\leq 7d_I(f,g)$ of Lemma~\\ref{Lem:FD leq Inter}. \nWhile that bound is obtained using an arbitrary selection, a better bound may be achievable using a particular optimal selection. \nIn addition, this may shed light on whether the bounds given between the bottleneck distance \nof the extended persistence diagrams and the two Reeb graph distances are tight. \nFinally, we will explore the applications of these distance measures to studying the stability of Reeb-like structures, such as Mapper and $\\alpha$-Reeb graphs. \n\n\n \n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{INTRODUCTION}\n\nIt is very interesting to study\nenergetic charged particle\ndiffusion in\nthe magnetic turbulence superposed on\nthe mean magnetic field\nin astrophysics\n (e.g.,\ncosmic ray physics, astrophysical\nplasmas, and\nspace weather research)\nand fusion plasma physics\n\\citep{Jokipii1966,\nSchlickeiser2002, MatthaeusEA2003,\nShalchiEA2005, ShalchiEA2006, Qin2007,\nHauffEA2008, Shalchi2009a,\nShalchi2010, QinEA2014}.\nThe interaction process between particles\nand turbulent magnetic fields is very\ncomplicated, and\nparticles and magnetic field lines\nhave the stochastic\nproperties. Therefore, one has to employ\nmethods of statistics to describe\nthe complicated transport of energetic particles\n\\citep{Earl1974, Earl1976,\nBeeckEA1986, Shalchi2011,\nLitvinenko2012a,\nLitvinenko2012b, ShalchiEA2013, HeEA2014,\nWangEA2017a, WangEA2017b, wq2018}\nand field line random walk\n\\citep{MatthaeusEA1995,\nShalchiAKourakis2007,\nShalchiAQin2010}.\nBecause the background\nmagnetic field $\\vec{B}_0$\nbreaks\nthe symmetry of the\nmagnetized plasma,\none has to distinguish\nparticle diffusion\nalong and across\nthe large-scale magnetic field.\nOne may\nonly consider the parallel\ndiffusion in many cases, when it is\nmuch greater than the perpendicular\none\n\\citep[see, e.g.,][]{Earl1974, Earl1976,\nBeeckEA1986, Shalchi2009b, Shalchi2011,\nLitvinenko2012a,\nLitvinenko2012b, ShalchiEA2013,\nHeEA2014, wq2018}.\n\nVarious analytical theories of\nparallel diffusion for energetic\ncharged particles have been\ndeveloped in the past.\nThe first attempt to resolve the parallel\ndiffusion was\nthe development of quasilinear theory\n\\citep{Jokipii1966} which corresponds\nto the first order\nperturbation theory.\nQuasilinear theory, however,\nis problematic and usually does\nnot agree with test-particle simulations\n\\citep{QinEA2002a, QinEA2002b,\nShalchi2009a}.\nWith the second order quasilinear theory\n(SOQLT) developed by \\citet{Shalchi2005},\nby using the Earl's formula\n\\citep{Earl1974}\nthe parallel\ndiffusion coefficient\ncan be evaluated\nmore accurately\n\\citep{ShalchiEA2009a,\nShalchiEA2009b, ReimerEA2016}.\nOn the other hand,\nby employing the assumptions\nand approximations\nused in NonLinear Guiding Center (NLGC)\ntheory \\citep{MatthaeusEA2003},\n\\citet{Qin2007} developed an extended\nparallel diffusion theory.\n\nThe mean square displacement\n$\\sigma^2=\\langle (\\Delta z)^2 \\rangle\n-\\langle (\\Delta z) \\rangle^2$\nis defined with the first and second order parallel\ndisplacement moments of the isotropic\ndistribution function\nof the energetic particles.\n$\\sigma^2$, describing the spread of\npossible particle orbits,\nis one of the fundamental quantities in the parallel transport.\nIn previous research\n\\citep[see, e.g.,][]{Shalchi2009a}\nthe parallel diffusion\ncoefficient $\\kappa_{zz}$\nis usually described by\nthe following formula\n(hereafter DCDV)\n\\begin{equation}\n\\kappa_{zz}=\\frac{1}{2}\\frac{d\\sigma^2}{dt}.\n\\label{kzz-sigma}\n\\end{equation}\n\nTo study energetic particle diffusion,\nthe background magnetic field $\\vec{B}_0$,\ne.g., mean solar wind magnetic field\nin the interplanetary space,\nis usually considered as a constant.\nHowever,\nit is clear that the mean solar wind\nmagnetic field\nis not constant in reality,\nespecially when particles are\nclose to the Sun.\nIt is found\nthat the\nspatially varying\nbackground solar wind magnetic fields lead\nto the adiabatic focusing effect of\ncharged energetic particle\ntransport and introduce\ncorrection to the particle\ndiffusion coefficients\n\\citep[see, e.g.,][]{Roelof1969,\nEarl1976, Kunstmann1979,BeeckEA1986,\nBieberEA1990,\nKota2000, SchlickeiserEA2008,\nShalchi2009b, Shalchi2011,\nLitvinenko2012a,\nLitvinenko2012b, ShalchiEA2013,\nWangEA2016, WangEA2017b, wq2018}.\nThe adiabatic focusing effect\ncauses\na convection term in\nthe energetic particles transport equation\nof the isotropic distribution function,\nso DCDV might be modified.\n\nIn order to address this question,\nin this paper we derive a new formula\nfor the parallel diffusion coefficient\nwith the parallel displacement variance\nby modifying DCDV with the\nadiabatic focusing\neffect.\nThe paper is organized as follows.\nIn Section \\ref{EQUATION OF ISOTROPIC\nDISTRIBUTION FUNCTION},\nfrom the Fokker-Planck equation\nwith adiabatic focusing effect\n we derive\nthe isotropic\ndistribution function equation\nand the telegraph equation.\nIn Section \\ref{THE PARALLEL DIFFUSION\nCOEFFICIENT REPRESENTED BY\nALONG-FIELD DISPLACEMENT MOMENTS},\nwe evaluate the first\nand second order moments of\nthe parallel displacement\nto derive the modified DCDV.\nIn Section \\ref{SPECIAL CASES},\nsome special cases of\nthe modified DCDV\nare explored.\nWe conclude and summarize\nour results in Section\n\\ref{SUMMARY AND CONCLUSION}.\n\n\\section{EQUATION OF ISOTROPIC\nDISTRIBUTION FUNCTION}\n\\label{EQUATION OF ISOTROPIC\nDISTRIBUTION FUNCTION}\n\nThe standard Fokker-Planck equation,\nwhich incorporates the effects\nof pitch-angle\nscattering and along-field\nadiabatic focusing,\nis given as\n\\citep{Roelof1969, Earl1981}\n\\begin{equation}\n\\frac{\\partial{f_0}}{\\partial{t}}\n+ v\\mu\n\\frac{\\partial{f_0}}\n{\\partial{z}}\n=\\frac{\\partial{}}\n{\\partial{\\mu}}\n\\left[D_{\\mu \\mu}(\\mu)\n\\frac{\\partial{f_0}}\n{\\partial{\\mu}}\\right]\n-\\frac{v}{2L}(1-\\mu^2)\n\\frac{\\partial{f_0 }}\n{\\partial{\\mu}},\n\\label{standard Fokker-Planck\nequation}\n\\end{equation}\nwhich is usually used\nin previous research\n\\citep{SchlickeiserEA2007,\nShalchi2011,\nLitvinenkoASchlickeiser2013,\nEffenbergerEA2014,\nMalkovSagdeev2015,\nWangEA2016}.\nHere $f_0$ is the isotropic\ndistribution function,\n$t$ is time, $z$\nis the distance\nalong the background\nmagnetic field,\n$\\mu=v_z \/v$ is the pitch-angle\ncosine with particle\nspeed $v$ and\nits z-component $v_z$,\n$D_{\\mu \\mu}(\\mu)$\nis the\npitch-angle diffusion\ncoefficient,\n$L(z)=-B_0 (z)\/ [dB_0 (z) \/ dz]$\nis the adiabatic focusing\ncharacteristic length of the\nlarge-scale magnetic field $B_0(z)$.\nIn this paper,\nfor simplification\nwe assume that\nthe adiabatic focusing\ncharacteristic length\nis a constant.\nThe source term and\nthe terms related to\nmomentum diffusion\nand so on are ignored\nin Equation\n(\\ref{standard Fokker-Planck\nequation}).\nThe more complete form of\nthe Fokker-Planck equation\ncan be found in\n\\citet{Schlickeiser2002}.\n\nBy introducing the linear\nphase space density\n$f(z,\\mu,t)=f_0(z,\\mu,t)\/B_0$,\nfrom Equation\n(\\ref{standard Fokker-Planck\nequation}),\nthe modified Fokker-Planck equation\nfor the\ndistribution function of energetic\ncharged particles can be obtained\n\\citep{Kunstmann1979,HeEA2014,\nWangEA2017b, wq2018}\n\\begin{equation}\n\\frac{\\partial{f}}{\\partial{t}}\n+ v\\mu\n\\frac{\\partial{f}}\n{\\partial{z}}=\n\\frac{\\partial{}}{\\partial{\\mu}}\n\\left[D_{\\mu \\mu}(\\mu)\n\\frac{\\partial{f}}\n{\\partial{\\mu}}-\\frac{v}{2L}\n(1-\\mu^2)f \\right].\n\\label{Fokker-Planck equation-0}\n\\end{equation}\nHere, $D_{\\mu\\mu}(\\mu)$ is time-independent. However,\nif the temporal characteristic\nof the pitch angle diffusion\nneed to be considered, $D_{\\mu\\mu}=D_{\\mu\\mu}(\\mu,t)$,\nEquation\n(\\ref{Fokker-Planck equation-0})\ncan be generalized as\n\\begin{equation}\n\\frac{\\partial{f}}{\\partial{t}}\n+ v\\mu\n\\frac{\\partial{f}}\n{\\partial{z}}=\n\\frac{\\partial{}}{\\partial{\\mu}}\n\\left[D_{\\mu \\mu}(\\mu,t)\n\\frac{\\partial{f}}\n{\\partial{\\mu}}-\\frac{v}{2L}\n(1-\\mu^2)f \\right].\n\\label{Fokker-Planck equation}\n\\end{equation}\nIn order for the pitch angle diffusion equation to hold,\nthe characteristic\ntime of the pitch angle\ndiffusion has to\nbe much less than\nthe characteristic\ntime of $D_{\\mu\\mu}(\\mu,t)$.\nEquation (\\ref{Fokker-Planck equation}) is the starting\npoint of the research in\nthis paper.\n\nWith strong pitch-angle scattering,\nthe gyro-tropic\ncosmic-ray phase space\ndensity $f(\\vec{x}, \\mu, t)$\ncan be split into\nthe dominant isotropic part\n$F(\\vec{x}, t)$\nand the subordinate\nanisotropic part\n$g(\\vec{x}, \\mu, t)$\n\\citep[see, e.g.,][]\n{SchlickeiserEA2007,\nSchlickeiserEA2008, HeEA2014,\nWangEA2017b, wq2018}\n\\begin{equation}\nf(\\vec{x}, \\mu, t)=F(\\vec{x}, t)\n+g(\\vec{x}, \\mu, t)\n\\end{equation}\nwith\n\\begin{equation}\nF(\\vec{x}, t)=\\frac{1}{2}\n\\int_{-1}^1 d\\mu\nf(\\vec{x}, \\mu, t)\n\\end{equation}\nand\n\\begin{equation}\n\\int_{-1}^1 d\\mu\ng(\\vec{x}, \\mu, t)=0 .\n\\end{equation}\n\n\\subsection{The variable\ncoefficient differential\nequation of the\nisotropic distribution function\n$F(z, t)$}\n\\label{The governing equation\nof the\nisotropic distribution function\nF}\n\nIn this subsection,\nby employing the method\nin \\citet{wq2018}\nwe derive\nthe formula of the variable\ncoefficient differential\nequation of the\nisotropic distribution function\n$F(z, t)$, and the derivation are\nsimilar to Equation (22) in\n\\citet{wq2018}.\n\nIntegration of Equation\n (\\ref{Fokker-Planck equation})\nover $\\mu$ gives\n\\begin{equation}\n\\frac{\\partial{F}}{\\partial{t}}\n+ \\frac{v}{2}\n\\frac{\\partial{}}\n{\\partial{z}}\\int_{-1}^{1}\n\\mu g d\\mu=0.\n\\label{Equation of F with g}\n\\end{equation}\nNext, integrating Equation\n(\\ref{Fokker-Planck equation})\nover $\\mu$ from $-1$ to $\\mu$,\nand using the regularity\n$D_{\\mu\\mu}(\\mu=\\pm 1)=0$\nwe can find\n\\begin{eqnarray}\n\\frac{\\partial{F}}\n{\\partial{t}}(\\mu+1)\n&&+ \\frac{\\partial{}}\n{\\partial{t}}\n\\int_{-1}^{\\mu}d\\nu g\n+\\frac{v(\\mu^2-1)}{2}\n\\frac{\\partial{F}}\n{\\partial{z}}\n+v\\frac{\\partial{}}\n{\\partial{z}}\n\\int_{-1}^{\\mu}\nd\\nu \\nu g\n=D_{\\mu\\mu}(\\mu,t)\n\\frac{\\partial{g}}\n{\\partial{\\mu}}\n-\\frac{v(1-\\mu^2)}{2L}F\n-\\frac{v(1-\\mu^2)}{2L}g.\n\\label{integrate\nfrom -1 to mu}\n\\end{eqnarray}\nBy subtracting Equation\n(\\ref{Equation of F with g})\nfrom\n(\\ref{integrate from -1 to mu})\nwe can find the\nfollowing equation\n\\begin{eqnarray}\n\\frac{\\partial{F}}\n{\\partial{t}}\\mu\n+ \\frac{\\partial{}}\n{\\partial{t}}\n\\int_{-1}^{\\mu}d\\nu g\n&&+\\frac{v(\\mu^2-1)}{2}\n\\frac{\\partial{F}}\n{\\partial{z}}\n+v\\frac{\\partial{}}\n{\\partial{z}}\n\\int_{-1}^{\\mu}\nd\\nu \\nu g\n-\\frac{v}{2}\n\\frac{\\partial{}}\n{\\partial{z}}\\int_{-1}^{1}\n\\mu g d\\mu \\nonumber\\\\\n&&=D_{\\mu\\mu}(\\mu,t)\n\\frac{\\partial{g}}\n{\\partial{\\mu}}\n-\\frac{v(1-\\mu^2)}{2L}F\n-\\frac{v(1-\\mu^2)}{2L}g.\n\\label{integrate from\n-1 to mu -2}\n\\end{eqnarray}\nAfter a straightforward\nalgebra, Equation\n(\\ref{integrate from -1 to mu -2})\nreduces to the following form\n\\begin{equation}\n\\frac{\\partial{g}}\n{\\partial{\\mu}}-\n\\frac{v(1-\\mu^2)g}\n{2L D_{\\mu \\mu}(\\mu,t)}\n+\\frac{v(1-\n\\mu^2)}{2D_{\\mu \\mu}(\\mu,t)}\n\\left(\\frac{\\partial{F}}\n{\\partial{z}}\n-\\frac{F}\n{L}\\right)=\\Phi(\\mu,t)\n\\label{Equation of g with phi}\n\\end{equation}\nwith\n\\begin{eqnarray}\n\\Phi(\\mu,t)=\n\\frac{1}{D_{\\mu\\mu}(\\mu,t)}\n\\Bigg[\\left(\\frac{\\partial{F}}\n{\\partial{t}}\\mu\n+\\frac{\\partial{}}{\\partial{t}}\n\\int_{-1}^{\\mu}gd\\nu\\right\n+\\frac{v}{2}\\frac{\\partial{}}\n{\\partial{z}}\n\\left(2\\int_{-1}^{\\mu}d\\nu \\nu g-\n\\int_{-1}^{1}d\\mu \\mu g\\right)\\Bigg].\n\\label{Phi}\n\\end{eqnarray}\nEquation\n(\\ref{Equation of g with phi})\ncan be rewritten as\n\\begin{equation}\n\\frac{\\partial{}}{\\partial{\\mu}}\n\\Bigg\\{\\Bigg[g(\\mu,t)\n-L\\left(\\frac{\\partial{F}}\n{\\partial{z}}\n-\\frac{F}{L} \\right)\n\\Bigg]e^{-M(\\mu,t)}\\Bigg\\}\n=e^{-M(\\mu,t)}\\Phi(\\mu,t)\n\\label{HS-like}\n\\end{equation}\nwith\n\\begin{equation}\nM(\\mu,t)=\\frac{v}{2L}\n\\int_{-1}^{\\mu} d\\nu\n\\frac{1-\\nu^2}{D_{\\nu \\nu}(\\nu,t)}.\n\\label{M(mu)}\n\\end{equation}\n\nIn additon, with the same procedure used\nin \\citet{wq2018},\nthe anisotropic\ndistribution function can be\nobtained as following\n\\begin{equation}\ng(\\mu,t)=L\\left(\\frac{\\partial{F}}\n{\\partial{z}}\n-\\frac{F}{L}\\right)\\left[1-\n\\frac{2e^{M(\\mu,t)}}{\\int_{-1}^{1}\nd\\mu e^{M(\\mu,t) }}\\right]\n+e^{M(\\mu,t)}\\left[R(\\mu,t)\n-\\frac{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t)}R(\\mu,t)}\n{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t) }}\\right]\n\\label{g}\n\\end{equation}\nwith\n\\begin{eqnarray}\nR(\\mu,t)&=&\\int_{-1}^{\\mu} d\\nu\ne^{-M(\\nu,t)}\\Phi(\\nu,t).\n\\label{R(mu)}\n\\end{eqnarray}\nEquations (\\ref{g}) and (\\ref{R(mu)})\ncontain\nthe effect from the term on the\nright hand side of\nEquation (\\ref{HS-like}).\nCombining Equations (\\ref{Phi}),\n(\\ref{g}) and (\\ref{R(mu)})\ngives the iterative function of\n$g(\\mu,t)$,\n\\begin{eqnarray}\ng(\\mu,t)&=&L\\left(\\frac{\\partial{F}}\n{\\partial{z}}\n-\\frac{F}{L}\\right)\\left[1-\n\\frac{2e^{M(\\mu,t)}}{\\int_{-1}^{1}\nd\\mu e^{M(\\mu,t) }}\\right]\n+e^{M(\\mu,t)}\\Bigg\\{\\int_{-1}^{\\mu}\nd\\nu\ne^{-M(\\nu,t)}\n\\frac{1}{D_{\\nu\\nu}(\\nu,t)}\n\\Bigg[\\left(\\frac{\\partial{F}}\n{\\partial{t}}\\nu\n+\\frac{\\partial{}}{\\partial{t}}\n\\int_{-1}^{\\nu}\ng(\\rho,t)d\\rho\\right)\n\\nonumber \\\\\n&&+\\frac{v}{2}\\frac{\\partial{}}\n{\\partial{z}}\n\\left(2\\int_{-1}^{\\nu}d\\rho \\rho\ng(\\rho,t)-\n\\int_{-1}^{1}d\\mu \\mu\ng(\\mu,t)\\right)\n\\Bigg]\\nonumber\\\\\n&&-\\frac{1}\n{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t) }}\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t)}\\int_{-1}^{\\mu} d\\nu\ne^{-M(\\nu,t)}\n\\frac{1}{D_{\\nu\\nu}(\\nu,t)}\n\\Bigg[\\left(\\frac{\\partial{F}}\n{\\partial{t}}\\nu\n+\\frac{\\partial{}}{\\partial{t}}\n\\int_{-1}^{\\nu}g(\\rho,t)\nd\\rho\\right)\\nonumber\\\\\n&&+\\frac{v}{2}\n\\frac{\\partial{}}\n{\\partial{z}}\n\\left(2\\int_{-1}^{\\mu}\nd\\nu \\nu g(\\nu,t)-\n\\int_{-1}^{1}d\\mu \\mu\ng(\\mu,t)\\right)\\Bigg]\\Bigg\\}.\n\\label{g2}\n\\end{eqnarray}\nBy iterating Equation\n(\\ref{g2}),\nwe obtain the formula of\n$g(\\mu,t)$ as\n\\begin{equation}\ng(\\mu,t)=\\sum_{m, n}\n\\epsilon_{m,n}(t)\n\\frac{\\partial^{m+n}{}}\n{\\partial {t^m}\n\\partial{z}^n} F\n\\label{expanded g}\n\\end{equation}\nwith the time dependent coefficients\n$\\epsilon_{m,n}(t)$.\nHere $m,n=0,1,2,3,\\cdots$.\nFurthermore,\nconsidering Equation\n(\\ref{Phi}), (\\ref{g})\nand Equation (\\ref{R(mu)}),\nwe can write\n$R(\\mu,t)$ as\n\\begin{equation}\nR(\\mu,t)=\\sum_{m, n}\n\\chi_{m,n}(t)\n\\frac{\\partial^{m+n}{}}\n{\\partial {t^m}\n\\partial{z}^n} F\n\\label{series of R}\n\\end{equation}\nwith the time dependent coefficients\n$\\chi_{m,n}(t)$ and\n$m,n=0,1,2,3,\\cdots$.\nThe coefficients\n$\\chi_{m,n}(t)$\nin the latter equation\nis related to\nthe coefficients\n$\\epsilon_{m,n}(t)$\nin Equation\n(\\ref{expanded g}).\n\nIn order to\nobtain the differential equation\nof isotropic distribution\nfunction,\nwe need to get the\nfollowing formula\nby using Equation\n(\\ref{Equation of F with g})\n\\begin{equation}\n\\int_{-1}^{1}d\\mu \\mu g(\\mu,t)\n=-2\\frac{\\int_{-1}^{1}d\\mu\\mu\ne^{M(\\mu,t)}}\n{\\int_{-1}^{1}d\\mu e^{M(\\mu,t)}}\n\\left(\\frac{\\partial{F}}\n{\\partial{z}}\n-\\frac{F}{L}\\right)L\n+\\int_{-1}^{1}\nd\\mu \\mu e^{M(\\mu,t)}\n\\Bigg[R(\\mu,t)\n-\\frac{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t)}R(\\mu,t)}\n{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t)}}\n\\Bigg].\n\\label{Integrating mu g over\nmu from -1 to 1}\n\\end{equation}\nBy inserting Equation\n(\\ref{Integrating mu g over\nmu from -1 to 1})\ninto Equation\n(\\ref{Equation of F with g})\none can obtain\n\\begin{eqnarray}\n\\frac{\\partial{F}}{\\partial{t}}\n&-&\\frac{\\partial{}}\n{\\partial{z}}\n\\left[vL\\frac{\\int_{-1}^{1}\nd\\mu \\mu e^{M(\\mu,t)}}\n{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t)}}\n\\left(\\frac{\\partial{F}}\n{\\partial{z}}-\n\\frac{F}{L}\\right)\\right]\n+\\Lambda(z,t)=0\n\\label{Accurate\ndiffusion equation}\n\\end{eqnarray}\nwith\n\\begin{equation}\n\\Lambda(z,t)=\n\\frac{v}{2}\\int_{-1}^{1}\nd\\mu \\mu e^{M(\\mu,t)}\n\\left[\\frac{\\partial{R(\\mu,t)}}\n{\\partial{z}}\n-\\frac{\\int_{-1}^{1}d\\mu\n\\frac{\\partial{R(\\mu,t)}}\n{\\partial{z}}\ne^{M(\\mu,t)}}\n{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t)}}\n\\right].\n\\label{Lambda}\n\\end{equation}\nBy inserting Equation\n(\\ref{series of R})\ninto Equation (\\ref{Lambda})\nwe can find\n\\begin{equation}\n\\Lambda(z,t)\n=\\sum_{m,n}\\eta_{m,n}(t)\n\\frac{\\partial^{m+n}}\n{\\partial {t^m}\n\\partial{z}^n} F\n\\label{Lambda-eta}\n\\end{equation}\nwith the coefficients\n$\\eta_{m,n}(t)$ and\n$m,n=0,1,2,3,\\cdots$.\nThus, Equation\n(\\ref{Accurate diffusion\nequation}) can be written as\n\\begin{eqnarray}\n\\frac{\\partial{F}}\n{\\partial{t}}\n-\\frac{\\partial{}}\n{\\partial{z}}&&\n\\left[vL\\frac{\\int_{-1}^{1}\nd\\mu \\mu e^{M(\\mu, t)}}\n{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t)}}\n\\left(\\frac{\\partial{F}}\n{\\partial{z}}-\n\\frac{F}{L}\\right)\\right]\n+\\sum_{m,n}\\eta_{m,n}(t)\n\\frac{\\partial^{m+n}}\n{\\partial {t^m}\n\\partial{z}^n} F=0.\n\\label{Accurate diffusion\nequation-2}\n\\end{eqnarray}\nNow we get\na variable coefficient\nlinear differential equation\nof the isotropic distribution\nfunction.\n\n\\subsection{The specific form of the variable coefficient differential equation}\n\nBy employing the method in\n\\citet{wq2018},\none can rewrite Equation\n(\\ref{Accurate diffusion\n\tequation-2})\nas\n\\begin{eqnarray}\n\\frac{\\partial{F}}\n{\\partial{t}}=\n&&-\\kappa_z(t)\n\\frac{\\partial{F}}{\\partial{z}}\n+\\kappa_{zz}(t)\n\\frac{\\partial^2{F}}{\\partial{z^2}}\n+\\left(\\kappa_{zzz}(t)\n\\frac{\\partial^3{F}}{\\partial{z^3}}\n+\\kappa_{zzzz}(t)\n\\frac{\\partial^4{F}}{\\partial{z^4}}\n+\\cdots\\right)\n+\\left( \\kappa_{tz}(t)\n\\frac{\\partial^2{F}}{\\partial{t}\n\\partial{z}}\n+ \\kappa_{ttz}(t)\n\\frac{\\partial^3{F}}\n{\\partial{t^2}\\partial{z}}\n+ \\kappa_{tttz}(t)\n\\frac{\\partial^4{F}}{\\partial{t^3}\n\\partial{z}}+\\cdots\\right)\n\\nonumber\\\\\n&&+\\left( \\kappa_{tzz}(t)\n\\frac{\\partial^3{F}}{\\partial{t}\n\\partial{z^2}}\n+ \\kappa_{ttzz}(t)\\frac{\\partial^4{F}}\n{\\partial{t^2}\\partial{z^2}}\n+ \\kappa_{tttzz}(t)\\frac{\\partial^5{F}}\n{\\partial{t^3}\\partial{z^2}}\n+\\cdots\\right)\n+\\cdots\\cdots,\n\\label{equation of F}\n\\end{eqnarray}\nhere, $\\kappa_z(t)$,\n$\\kappa_{zz}(t)$,\n$\\kappa_{zzz}(t)$,\n$\\kappa_{zzzz}(t)$,\n$\\cdots\\cdots$,\n$\\kappa_{tz}(t)$,\n$\\kappa_{ttz}(t)$, $\\kappa_{tttz}(t)$,\n$\\cdots\\cdots$,\n$\\kappa_{tzz}(t)$, $\\kappa_{ttzz}(t)$,\n$\\kappa_{tttzz}(t)$,\n$\\cdots\\cdots$\nare all time dependent.\nThis equation determines the transport regimes by the specific\nforms of its coefficient.\nFor example, as shown in\n\\citet{wq2018},\n$\\kappa_z(t)$\nis the coefficient of the convection\nterm\n\\begin{equation}\n\\kappa_z(t)\n=v\\frac{\\int_{-1}^{1}d\\mu\n\\mu e^{M(\\mu,t)}}{\\int_{-1}^{1}\nd\\mu e^{M(\\mu,t)}},\n\\label{kz}\n\\end{equation}\nand $\\kappa_{zz}(t)$\nis the coefficient of the parallel\ndiffusion term\n\\begin{eqnarray}\n\\kappa_{zz}(t)&=&\nvL\\frac{\\int_{-1}^{1}d\\mu\n\\mu e^{M(\\mu,t)}}{\\int_{-1}^{1}\nd\\mu e^{M(\\mu,t)}}\n+\\frac{v^2}{2}\n\\Bigg\\{\n\\int_{-1}^{1}d\\mu\n\\mu e^{M(\\mu,t)}\n\\int_{-1}^{\\mu}d\\nu\n\\frac{e^{-M(\\nu,t)}}\n{D_{\\nu \\nu}(\\nu,t)}\n\\left[\\int_{-1}^{\\nu}d\\rho\\rho\n\\left(1-\\frac{2e^{M(\\rho,t)}}\n{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t)}}\\right)\n+\\int_{-1}^{1} d\\mu \\mu\n\\frac{e^{M(\\mu,t)}}\n{\\int_{-1}^{1} d\\mu\ne^{M(\\mu,t)}}\\right]\n\\nonumber\\\\\n&&-\\frac{\\int_{-1}^{1}\nd\\mu \\mu\ne^{M(\\mu,t)}}{\\int_{-1}^{1}\nd\\mu e^{M(\\mu,t)}}\n\\int_{-1}^{1}\nd\\mu e^{M(\\mu,t)}\n\\int_{-1}^{\\mu}\nd\\nu \\frac{e^{-M(\\nu,t)}}\n{D_{\\nu \\nu}(\\nu,t)}\n\\left[\\int_{-1}^{\\nu}\nd\\rho\\rho\n\\left(1-\\frac{2e^{M(\\rho,t)}}\n{\\int_{-1}^{1}d\\mu\ne^{M(\\mu,t)}}\\right)\n+\\int_{-1}^{1} d\\mu \\mu\n\\frac{e^{M(\\mu,t)}}\n{\\int_{-1}^{1} d\\mu\ne^{M(\\mu,t)}}\\right]\n\\Bigg\\}.\n \\label{A1-iso}\n\\end{eqnarray}\nWe can similarly obtain other\ntime-dependent coefficients in Equation\n(\\ref{equation of F}).\n\n\\subsection{Telegraph equation of the\nisotropic distribution function\n$F(z, t)$}\n\\label{Telegraph equation of the\nisotropic distribution function}\n\nThe diffusion convection equation\nis frequently used in previous research\n\\citep[see, e.g.,][]{Shalchi2009a}.\nBut the accuracy of the diffusion\napproximation is limited\nbecause the signal speed is\ninfinite in the diffusion limit.\nRecently, an improved description\nof the\nenergetic particle transport\nis provided by the telegraph equation,\ncharacterized by a finite\nsignal propagation speed\n\\citep{Kota1994,\nPaulsBurger1994,\nSchwadronEA2010, PorraMasoliver1997,\nZankEA2000, FedorovShakhov2003,\nKaghashviliEA2004,\nLitvinenkoASchlickeiser2013,\nEffenbergerEA2014, LitvinenkoEA2015,\nLitvinenkoNoble2016}.\nUsing the iterative method\nin \\citet{LitvinenkoASchlickeiser2013},\none can obtain the\ntelegraph equation from\nEquation (\\ref{equation of F})\n\\begin{eqnarray}\n\\frac{\\partial{F}}{\\partial{t}}\n+\\tau(t)\\frac{\\partial^2{F}}{\\partial{t^2}\n}+\\kappa_z(t)\\frac{\\partial{F}}\n{\\partial{z}}=\\kappa_{zz}(t)\n\\frac{\\partial^2{F}}{\\partial{z^2}}\n\\label{telegraph equation}\n\\end{eqnarray}\nwith\n\\begin{eqnarray}\n\\tau(t)=-\\frac{\\kappa_{tzz}(t)}\n{\\kappa_{zz}(t)}.\n\\label{tau}\n\\end{eqnarray}\n\nEquation (\\ref{equation of F})\ncan be rewritten as\n\\begin{eqnarray}\n\\frac{\\partial{F}}{\\partial{z}}\n=&&\\frac{1}{\\kappa_z(t)}\n\\Bigg[-\\frac{\\partial{F}}\n{\\partial{t}}\n+\\kappa_{zz}(t)\n\\frac{\\partial^2{F}}{\\partial{z^2}}\n+\\left(\\kappa_{zzz}(t)\n\\frac{\\partial^3{F}}{\\partial{z^3}}\n+\\kappa_{zzzz}(t)\n\\frac{\\partial^4{F}}{\\partial{z^4}}\n+\\cdots\\right)\n+\\left( \\kappa_{tz}(t)\n\\frac{\\partial^2{F}}{\\partial{t}\n\t\\partial{z}}\n+ \\kappa_{ttz}(t)\n\\frac{\\partial^3{F}}\n{\\partial{t^2}\\partial{z}}\n+ \\kappa_{tttz}(t)\n\\frac{\\partial^4{F}}{\\partial{t^3}\n\t\\partial{z}}+\\cdots\\right)\n\\nonumber\\\\\n&&+\\left( \\kappa_{tzz}(t)\n\\frac{\\partial^3{F}}{\\partial{t}\n\t\\partial{z^2}}\n+ \\kappa_{ttzz}(t)\\frac{\\partial^4{F}}\n{\\partial{t^2}\\partial{z^2}}\n+ \\kappa_{tttzz}(t)\\frac{\\partial^5{F}}\n{\\partial{t^3}\\partial{z^2}}\n+\\cdots\\right)\n+\\cdots\\cdots\\Bigg].\n\\label{convection term of equation of F}\n\\end{eqnarray}\nInserting Equation\n(\\ref{convection term of equation of F}) into the convection term\n$\\kappa_{tz}\n\\partial^2{F}\/\\partial{t}\\partial{z}$\nin Equation (\\ref{equation of F})\nand neglecting the higher order derivative terms,\nwe can obtain\nanother version of\nthe telegraph equation,\n\\begin{eqnarray}\n\\frac{\\partial{F}}{\\partial{t}}\n+\\tau'(t)\\frac{\\partial^2{F}}{\\partial{t^2}}\n+\\kappa_z(t)\\frac{\\partial{F}}{\\partial{z}}\n=\\kappa_{zz}(t)\\frac{\\partial^2{F}}\n{\\partial{z^2}}\n\\label{new telegraph equation}\n\\end{eqnarray}\nwith\n\\begin{eqnarray}\n\\tau'(t)=\\frac{\\kappa_{tz}\n(t)}{\\kappa_{z}(t)}.\n\\label{tau}\n\\end{eqnarray}\nWe find that the only\ndifference between Equations (\\ref{new telegraph equation})\nand (\\ref{telegraph equation})\nis caused by the fact that\n$\\tau(t)$ is not equal to $\\tau'(t)$.\n\n\\section{THE MODIFIED FORMULA OF DIFFUSION COEFFICIENT\nWITH DISPLACEMENT VARIANC}\n\\label{THE PARALLEL DIFFUSION\nCOEFFICIENT REPRESENTED\nBY ALONG-FIELD DISPLACEMENT MOMENTS}\nIn the following, we derive\nthe modified formula\nof the parallel\ndiffusion coefficient\nwith the parallel displacement variance\ninduced by the adiabatic focusing\n(hereafter MDCDV).\n\n\\subsection{MDCDV\nfor the variable coefficient\ndifferential equation}\n\\label{The parallel diffusion\ncoefficient of Equation}\n\nBy multiplying Equation\n(\\ref{equation of F})\nwith\n$\\langle (\\Delta z) \\rangle$ and\nintegrating the result, one can find\n\\begin{eqnarray}\n\\frac{d}{dt}\\langle (\\Delta z)\n\\rangle&=&\\int_{-\\infty}^\\infty\ndz (\\Delta z)\\frac{\\partial{F}}\n{\\partial{t}}\\nonumber\\\\\n&=&\\int_{-\\infty}^\\infty dz\n(\\Delta z)\\Bigg[\\left(-\\kappa_z(t)\n\\frac{\\partial{F}}{\\partial{z}}\n+\\kappa_{zz}(t)\n\\frac{\\partial^2{F}}{\\partial{z^2}}\n+\\kappa_{zzz}(t)\n\\frac{\\partial^3{F}}{\\partial{z^3}}\n+\\kappa_{zzzz}(t)\n\\frac{\\partial^4{F}}{\\partial{z^4}}\n+\\cdots+\\kappa_{z\\cdots z}(t)\n\\frac{\\partial^n{F}}{\\partial{z^n}}\n+\\cdots\\right)\\nonumber\\\\\n&&+\\left( \\kappa_{tz}(t)\n\\frac{\\partial^2{F}}\n{\\partial{t}\\partial{z}}\n+ \\kappa_{ttz}(t)\n\\frac{\\partial^3{F}}\n{\\partial{t^2}\\partial{z}}\n+ \\kappa_{tttz}(t)\n\\frac{\\partial^4{F}}\n{\\partial{t^3}\\partial{z}}\n+\\cdots\n+\\kappa_{t\\cdots tz}(t)\n\\frac{\\partial^{n+1}{F}}\n{\\partial{t^n}\\partial{z}}\n+\\cdots\n\\right)\\nonumber\\\\\n&&+\\left( \\kappa_{tzz}(t)\n\\frac{\\partial^3{F}}\n{\\partial{t}\\partial{z^2}}\n+ \\kappa_{ttzz}(t)\n\\frac{\\partial^4{F}}\n{\\partial{t^2}\\partial{z^2}}+\n+\\cdots+\\kappa_{ttzz}(t)\n\\frac{\\partial^4{F}}\n{\\partial{t^2}\\partial{z^2}}\n+ \\kappa_{t\\cdots tzz}(t)\n\\frac{\\partial^{n+2}{F}}\n{\\partial{t^n}\\partial{z^2}}\n+\\cdots\\right)\n+\\cdots\\cdots\\Bigg].\n\\label{1-general}\n\\end{eqnarray}\nTo proceed,\nwe have to obtain the integrals\non the right hand side of\nthe latter equation\nwith the following conditions\n\\begin{eqnarray}\n&&F(z=\\pm\\infty)=0,\\\\\n&&\\frac{\\partial^n{F}}\n{\\partial{z^n}}(z=\\pm\\infty)=0\n\\hspace{0.2cm} n=1,2,3,\\cdots.\n\\end{eqnarray}\nBy performing integration\nin parts one can obtain\n\\begin{eqnarray}\n&&\\int_{-\\infty}^\\infty\ndz (\\Delta z)\\left(-\\kappa_z(t)\n\\frac{\\partial{F}}{\\partial{z}}\\right)\n=-\\kappa_z(t)\\int_{-\\infty}^\\infty\ndz (\\Delta z)\\frac{\\partial{F}}\n{\\partial{z}}=\\kappa_z(t),\n\\label{1-z}\\\\\n&&\\int_{-\\infty}^\\infty\ndz (\\Delta z)\\left(\\kappa_{zz}(t)\n\\frac{\\partial^2{F}}{\\partial{z^2}}\\right)\n=\\kappa_{zz}(t)\\int_{-\\infty}^\\infty\ndz (\\Delta z)\\frac{\\partial^2{F}}\n{\\partial{z^2}}=0.\n\\label{1-zz}\n\\end{eqnarray}\nWe can also find that the integrals\nare equal to zero if the order of spatial\nderivative of the integrands is\nhigher than one,\n\\begin{eqnarray}\n\\int_{-\\infty}^\\infty dz (\\Delta z)\n\\left(\\kappa_{z\\cdots z}(t)\n\\frac{\\partial^n{F}}{\\partial{z^n}}\\right)\n=\\kappa_{z\\cdots z}(t)\\int_{-\\infty}^\\infty\ndz (\\Delta z)\n\\frac{\\partial^n{F}}{\\partial{z^n}}\n=0 \\hspace{0.5cm}n=2,3,\\cdots.\n\\label{1-zn}\n\\end{eqnarray}\nIt is also easy to find the\nfollowing result by integration\nby parts\n\\begin{eqnarray}\n\\int_{-\\infty}^\\infty dz\n(\\Delta z)\\left(\\kappa_{z\\cdots z}(t)\n\\frac{\\partial^{n+m}{F}}{\\partial{t^m}\n\\partial{z^n}}\\right)=\\kappa_{z\\cdots z}(t)\n\\frac{d^m}{dt^m}\\int_{-\\infty}^\\infty\ndz (\\Delta z)\n\\frac{\\partial^n{F}}{\\partial{z^n}}\n=0 \\hspace{0.5cm}n=1,2,3,\\cdots,\n\\hspace{0.2cm}m=1,2,3,\\cdots.\n\\label{1-tmzn}\n\\end{eqnarray}\nInserting Equations (\\ref{1-z})\n- (\\ref{1-tmzn}) into Equation\n(\\ref{1-general}), one can obtain\nthe following formula\n\\begin{equation}\n\\frac{d}{dt}\\langle (\\Delta z)\n\\rangle=\\kappa_z(t).\n\\label{1-result}\n\\end{equation}\n\nFurthermore, analogous to\nEquation (\\ref{1-general}),\nthe formula of second order\nmoment of the parallel\ndisplacement\ncan be written as\n\\begin{eqnarray}\n\\frac{d}{dt}\\langle (\\Delta z)^2\n\\rangle&=&\\int_{-\\infty}^\\infty\ndz (\\Delta z)^2\n\\frac{\\partial{F}}{\\partial{t}}\n\\nonumber\\\\\n&=&\\int_{-\\infty}^\\infty dz\n(\\Delta z)^2\n\\Bigg[\\left(-\\kappa_z(t)\n\\frac{\\partial{F}}{\\partial{z}}\n+\\kappa_{zz}(t)\n\\frac{\\partial^2{F}}{\\partial{z^2}}\n+\\kappa_{zzz}(t)\n\\frac{\\partial^3{F}}{\\partial{z^3}}\n+\\kappa_{zzzz}(t)\n\\frac{\\partial^4{F}}{\\partial{z^4}}\n+\\cdots+\\kappa_{z\\cdots z}(t)\n\\frac{\\partial^n{F}}{\\partial{z^n}}\n+\\cdots\\right)\\nonumber\\\\\n&{}&+\\left( \\kappa_{tz}\n\\frac{\\partial^2{F}}\n{\\partial{t}\\partial{z}}\n+ \\kappa_{ttz}(t)\n\\frac{\\partial^3{F}}\n{\\partial{t^2}\\partial{z}}\n+ \\kappa_{tttz}(t)\n\\frac{\\partial^4{F}}\n{\\partial{t^3}\\partial{z}}\n+\\cdots\n+\\kappa_{t\\cdots tz}(t)\n\\frac{\\partial^{n+1}{F}}\n{\\partial{t^n}\\partial{z}}\n+\\cdots\\right)\\nonumber\\\\\n&&+\\left( \\kappa_{tzz}(t)\n\\frac{\\partial^3{F}}\n{\\partial{t}\\partial{z^2}}\n+ \\kappa_{ttzz}(t)\\frac{\\partial^4{F}}\n{\\partial{t^2}\\partial{z^2}}\n+ \\kappa_{tttzz}(t)\\frac{\\partial^5{F}}\n{\\partial{t^3}\\partial{z^2}}\n+\\cdots+\\kappa_{t\\cdots tzz}(t)\n\\frac{\\partial^{n+2}{F}}\n{\\partial{t^n}\\partial{z^2}}\n+\\cdots\\right)\n+\\cdots\\cdots\\Bigg].\n\\label{2-general}\n\\end{eqnarray}\nThe integrals on the right\nhand side of the latter\nequation\ncan be easily found\nby part integral\n\\begin{eqnarray}\n&&\\int_{-\\infty}^\\infty dz\n(\\Delta z)^2\\left(-\\kappa_z(t)\n\\frac{\\partial{F}}{\\partial{z}}\n\\right)\n=-\\kappa_z(t)\\int_{-\\infty}^\\infty\ndz (\\Delta z)^2\\frac{\\partial{F}}\n{\\partial{z}}\n=2\\kappa_z (t)\\langle (\\Delta z)\n\\rangle\n\\label{2-z}\\\\\n&&\\int_{-\\infty}^\\infty dz\n(\\Delta z)^2\\left(\\kappa_{zz}(t)\n\\frac{\\partial^2{F}}{\\partial{z^2}}\n\\right)\n=\\kappa_{zz}(t)\\int_{-\\infty}^\\infty\ndz (\\Delta z)^2\\frac{\\partial^2{F}}\n{\\partial{z^2}}\n=2\\kappa_{zz}(t)\n\\label{2-zz}\\\\\n&&\\int_{-\\infty}^\\infty dz\n(\\Delta z)^2\\left(\\kappa_{z\\cdots z}\n(t)\\frac{\\partial^n{F}}\n{\\partial{z^n}}\\right)\n=\\kappa_{z\\cdots z}(t)\n\\int_{-\\infty}^\\infty dz\n(\\Delta z)^2\\frac{\\partial^n{F}}\n{\\partial{z^n}}\n=0 \\hspace{0.5cm}n=3, 4, 5,\n\\cdots\\\\\n&&\\int_{-\\infty}^\\infty dz\n(\\Delta z)^2\n\\left(\\kappa_{t\\cdots tz}(t)\n\\frac{\\partial^{n+1}{F}}\n{\\partial{t^n}\\partial{z}}\\right)\n=\\kappa_{t\\cdots tz}(t)\n\\frac{d^n}{dt^n}\n\\int_{-\\infty}^\\infty dz\n(\\Delta z)^2\n\\frac{\\partial{F}}{\\partial{z}}\n=-2\\kappa_{t\\cdots tz}(t)\n\\frac{d^n}{dt^n}\\langle\n(\\Delta z) \\rangle\n\\hspace{0.5cm}n=1,2,3,\n\\cdots.\n\\label{1-tzn}\\\\\n&&\\int_{-\\infty}^\\infty dz\n(\\Delta z)^2\n\\left(\\kappa_{t\\cdots tz\\cdots z}(t)\n\\frac{\\partial^{n+m}{F}}\n{\\partial{t^n}\\partial{z^m}}\\right)\n=\\kappa_{t\\cdots tz\\cdots z}(t)\n\\frac{d^n}{dt^n}\n\\int_{-\\infty}^\\infty\ndz (\\Delta z)^2\n\\frac{\\partial{F^m}}\n{\\partial{z^m}}=0\n\\hspace{0.5cm}n=1,2,3,\n\\cdots, \\hspace{0.2cm}m=2,3,\n\\cdots.\n\\label{2-tzn}\n\\end{eqnarray}\nWith Equations\n(\\ref{2-z})-(\\ref{2-tzn}),\nEquation (\\ref{2-general})\nbecomes\n\\begin{eqnarray}\n\\frac{d}{dt}\\langle (\\Delta z)^2\n\\rangle =2\\kappa_{zz}(t)+2\\left[\\kappa_z(t)\n\\langle (\\Delta z) \\rangle\n-\\kappa_{tz}(t)\\frac{d}{dt}\n\\langle (\\Delta z) \\rangle\n-\\kappa_{ttz}(t)\\frac{d^2}{dt^2}\n\\langle (\\Delta z) \\rangle\n-\\kappa_{tttz}(t)\\frac{d^3}{dt^3}\n\\langle (\\Delta z) \\rangle\n-\\cdots\\cdots\\right].\n\\label{2-result}\n\\end{eqnarray}\nThe latter equation can be rewritten as\n\\begin{eqnarray}\n\\kappa_{zz}(t)=\\frac{1}{2}\n\\frac{d}{dt}\\langle\n(\\Delta z)^2 \\rangle\n+C(t)\n\\end{eqnarray}\nwith\n\\begin{eqnarray}\nC(t)=-\\kappa_z(t)\\langle (\\Delta z)\\rangle\n+\\kappa_{tz}(t)\\frac{d}{dt}\n\\langle (\\Delta z) \\rangle\n+\\kappa_{ttz}(t)\\frac{d^2}{dt^2}\n\\langle (\\Delta z) \\rangle\n+\\kappa_{tttz}(t)\\frac{d^3}{dt^3}\n\\langle (\\Delta z) \\rangle\n+\\cdots\\cdots.\n\\end{eqnarray}\nHere, $C(t)$\nis related to convection\nof energetic particles.\n\nBy employing formula $\\sigma^2=\\langle\n(\\Delta z)^2 \\rangle\n-\\langle (\\Delta z)\n\\rangle^2$ and Equation\n(\\ref{1-result})\n we can obtain\n\\begin{eqnarray}\n\\kappa_{zz}(t)=\\frac{1}{2}\n\\frac{d\\sigma^2}{dt}\n+\\left[\\kappa_{tz}(t)\\frac{d}{dt}\n\\langle (\\Delta z) \\rangle\n+\\kappa_{ttz}(t)\\frac{d^2}{dt^2}\n\\langle (\\Delta z) \\rangle\n+\\kappa_{tttz}(t)\\frac{d^3}{dt^3}\n\\langle (\\Delta z) \\rangle\n+\\cdots\\cdots\\right].\n\\label{center formula-1}\n\\end{eqnarray}\nInserting Equation (\\ref{1-result})\ninto the latter equation one\ncan find\n\\begin{eqnarray}\n\\kappa_{zz}(t)=\\frac{1}{2}\n\\frac{d\\sigma^2}{dt}\n+\\left[\\kappa_{tz}(t)\n\\kappa_z(t)\n+\\kappa_{ttz}(t)\\frac{d}{dt}\n\\kappa_z(t)\n+\\kappa_{tttz}(t)\\frac{d^2}{dt^2}\n\\kappa_z(t)+\\cdots\\cdots\\right].\n\\label{center formula-2}\n\\end{eqnarray}\nThe latter equation can\nbe rewritten as\n\\begin{eqnarray}\n\\frac{d\\sigma^2}{dt}\n=2\\kappa_{zz}(t)\n-\\left[2\\kappa_{tz}(t)\n\\kappa_z(t)\n+2\\kappa_{ttz}(t)\\frac{d}{dt}\n\\kappa_z(t)+2\\kappa_{tttz}(t)\n\\frac{d^2}{dt^2}\\kappa_z(t)\n+\\cdots\\cdots\\right].\n\\label{kzz-0}\n\\end{eqnarray}\n\nEquation (\\ref{center formula-1})\nor (\\ref{kzz-0}) is the modified DCDV\n(MDCDV).\nIn comparison with the well-known\nEquation (\\ref{kzz-sigma})\n(represented by DCDV),\none can find that MDCDV\nis a more general formula.\nThat is, DCDV\nis the special case of MDCDV\nwhen the cross terms with the first\norder spatial derivative or\nthe convection term are\nequal to zero.\n\nBecause the coefficients\n$\\kappa_{z}(t)$, $\\kappa_{zz}(t)$,\n$\\kappa_{tz}(t)$, $\\cdots\\cdots$\nare all time dependent,\nthe relation $\\sigma^2\\propto t$\ndoes not hold for MDCDV.\nTherefore,\nthe transport\ndescribed by MDCDV\nis not in diffusion regime.\nSince Equation (\\ref{kzz-0}) is\nderived from Equation\n(\\ref{equation of F}),\nthe transport described\nby Equation (\\ref{equation of F})\nis also not in diffusion regime.\n\n\\subsection{MDCDV\nfor the telegraph equation}\n\\label{The parallel diffusion\ncoefficient of the telegraph equation}\n\nFor the telegraph equation\nderived in\n\\citet{LitvinenkoASchlickeiser2013}\nthe corresponding first\nand second order moments of the\nparallel displacement\ncan be written as, respectively\n\\begin{eqnarray}\n&&\\frac{d}{dt}\\langle (\\Delta z)\n\\rangle=\\kappa_z(t)-\\tau(t)\n\\frac{d^2{}}{d{t^2}}\\langle\n(\\Delta z) \\rangle,\n\\label{1-result for\ntelegraph equation}\\\\\n&&\\frac{d}{dt}\\langle\n(\\Delta z)^2 \\rangle\n=2\\kappa_z(t) \\langle (\\Delta z)\n\\rangle+2\\kappa_{zz}(t)\n-\\tau(t)\\frac{d^2}{dt^2}\\langle\n(\\Delta z)^2 \\rangle.\n\\label{2-result for\ntelegraph equation}\n\\end{eqnarray}\nBy combining the latter equations\none can derive with $\\sigma^2=\\langle\n(\\Delta z)^2 \\rangle -\\langle\n(\\Delta z) \\rangle^2$\n\\begin{eqnarray}\n\\kappa_{zz}(t)&=&\\frac{1}{2}\n\\frac{d\\sigma^2}{dt}\n+\\tau(t)\\Bigg[\\frac{1}{2}\n\\frac{d^2}{dt^2}\\langle (\\Delta z)^2\n\\rangle-\\langle (\\Delta z)\n\\rangle\\frac{d^2}{dt^2}\\langle\n(\\Delta z) \\rangle\\Bigg].\n\\end{eqnarray}\nInserting Equation (\\ref{tau})\ninto the latter equation,\nwe can obtain\n\\begin{eqnarray}\n\\kappa_{zz}(t)\n=\\frac{1}{2}\\frac{d\\sigma^2}{dt}\n-\\frac{\\kappa_{tzz}(t)}\n{\\kappa_{zz}(t)}\\Bigg[\\frac{d^2}{dt^2}\n\\langle (\\Delta z)^2 \\rangle-\\langle\n(\\Delta z) \\rangle\\frac{d^2}{dt^2}\n\\langle (\\Delta z) \\rangle\\Bigg].\n\\label{kzz for telegraph equation}\n\\end{eqnarray}\n\nBy subtracting Equation\n(\\ref{kzz for telegraph equation})\nfrom Equation (\\ref{center formula-2})\nwe can obtain the error formula\n\\begin{eqnarray}\nE(t)\n=\\left[\\kappa_{tz}(t)\\kappa_z(t)\n+\\kappa_{ttz}(t)\\frac{d}\n{dt}\\kappa_z(t)\n+\\kappa_{tttz}(t)\\frac{d^2}{dt^2}\n\\kappa_z(t)+\\cdots\\cdots\\right]\n+\\frac{\\kappa_{tzz}(t)}\n{\\kappa_{zz}(t)}\\Bigg[\\frac{d^2}{dt^2}\n\\langle (\\Delta z)^2 \\rangle\n-\\langle (\\Delta z)\n\\rangle\\frac{d^2}{dt^2}\\langle\n(\\Delta z) \\rangle\\Bigg],\n\\label{difference}\n\\end{eqnarray}\nwhich comes from the iteration operation\nused in the derivation\nof telegraph equation.\nIn the future, this formula\nwill be explored in detail.\n\nIn addition, Equation\n(\\ref{1-result for telegraph equation})\ncan be rewritten as\n\\begin{eqnarray}\n\\tau(t)\\frac{d^2{}}{d{t^2}}\n\\langle (\\Delta z) \\rangle\n+\\frac{d}{dt}\\langle (\\Delta z)\n\\rangle=\\kappa_z(t).\n\\label{11-differential equation}\n\\end{eqnarray}\nIf the coefficients $\\tau(t)$\nand $\\kappa_z(t)$ can be determined,\nthe latter equation might\nbe solved.\nSimilarily, Equation\n(\\ref{2-result for\ntelegraph equation})\ncan be rewritten as\n\\begin{eqnarray}\n\\tau(t)\\frac{d^2}{dt^2}\n\\langle (\\Delta z)^2\n\\rangle+\\frac{d}{dt}\\langle\n(\\Delta z)^2 \\rangle =2\\kappa_z(t)\n\\langle (\\Delta z) \\rangle\n+2\\kappa_{zz}(t)\n\\label{22-differential equation}\n\\end{eqnarray}\nCombining the solution of\nEquation\n(\\ref{11-differential equation})\nand the coefficients\n$\\tau(t)$ and $\\kappa_z(t)$\nwith the latter equation,\nthe formula of\n$\\langle (\\Delta z)^2 \\rangle$\nmight be got.\nSome solutions\nof Equations (\\ref{11-differential equation})\nand (\\ref{22-differential equation})\nwill be investigated\nin the future.\n\n\\subsection{MDCDV for telegraph\nEquation (\\ref{new telegraph equation})}\n\\label{The parallel diffusion\ncoefficient of the new\ntelegraph equation}\n\nSimilarily, for telegraph\nEquation (\\ref{new telegraph equation}) the first and\nsecond moments can be got\nas follows\n\\begin{eqnarray}\n&&\\frac{d}{dt}\\langle (\\Delta z) \\rangle\n=\\kappa_z(t)-\\tau'(t)\\frac{d^2{}}\n{d{t^2}}\\langle (\\Delta z) \\rangle,\\\\\n&&\\frac{d}{dt}\\langle (\\Delta z)^2\n\\rangle =2\\kappa_z(t) \\langle (\\Delta z)\n\\rangle+2\\kappa_{zz}(t)\n-\\tau'(t)\\frac{d^2}{dt^2}\n\\langle (\\Delta z)^2 \\rangle.\n\\end{eqnarray}\nThe corresponding MDCDV becomes\n\\begin{eqnarray}\n\\kappa_{zz}(t)\n=\\frac{1}{2}\\frac{d\\sigma^2}{dt}\n+\\frac{\\kappa_{tz}(t)}{\\kappa_{z}(t)}\n\\Bigg[\\frac{d^2}{dt^2}\\langle (\\Delta z)^2\n\\rangle-\\langle (\\Delta z)\n\\rangle\\frac{d^2}{dt^2}\\langle\n(\\Delta z) \\rangle\\Bigg].\n\\end{eqnarray}\n\n\\section{SPECIAL CASES}\n\\label{SPECIAL CASES}\n\nIn the following, we investigate\nthe special cases of MDCDV\n(see Equation (\\ref{center formula-1})),\nwhich\ncorrespond to the\nsimplified versions of Equation\n(\\ref{equation of F}).\n\n\\subsection{Case 1:\nconstant coefficient\ndifferential equation}\n\\label{For the case D=D(mu)}\nIf the temporal effect of\nthe pitch-angle diffusion\ncoefficient $D_{\\mu\\mu}$\ncan be ignored,\ni.e.,\n$D_{\\mu\\mu}=D_{\\mu\\mu}(\\mu)$,\nall the coefficients in\nEquation (\\ref{equation of F})\nare constant.\nTherefore, Equation\n(\\ref{equation of F})\nis simplified as\n\\begin{eqnarray}\n\\frac{\\partial{F}}{\\partial{t}}=\n&&\\left(-\\kappa_z\\frac{\\partial{F}}\n{\\partial{z}}+\\kappa_{zz}\n\\frac{\\partial^2{F}}{\\partial{z^2}}\n+\\kappa_{zzz}\n\\frac{\\partial^3{F}}{\\partial{z^3}}\n+\\kappa_{zzzz}\n\\frac{\\partial^4{F}}{\\partial{z^4}}\n+\\cdots\\right)\n+\\left( \\kappa_{tz}\n\\frac{\\partial^2{F}}{\\partial{t}\n\\partial{z}}\n+ \\kappa_{ttz}\\frac{\\partial^3{F}}\n{\\partial{t^2}\\partial{z}}\n+ \\kappa_{tttz}\\frac{\\partial^4{F}}\n{\\partial{t^3}\\partial{z}}\n+\\cdots\\right)\\nonumber\\\\\n&&+\\left( \\kappa_{tzz}\n\\frac{\\partial^3{F}}{\\partial{t}\n\\partial{z^2}}\n+ \\kappa_{ttzz}\\frac{\\partial^4{F}}\n{\\partial{t^2}\\partial{z^2}}\n+ \\kappa_{tttzz}\\frac{\\partial^5{F}}\n{\\partial{t^3}\\partial{z^2}}\n+\\cdots\\right)\n+\\cdots\\cdots.\n\\label{equation of F with\nconstant coefficient}\n\\end{eqnarray}\nThis is a constant coefficient\nlinear differential equation\nof the isotropic distribution\nfunction.\nThe first\nand second order moments of\nthe parallel displacement can\nbe found as\n\\begin{eqnarray}\n&&\\frac{d}{dt}\\langle (\\Delta z)\n\\rangle=\\kappa_z=constant,\n\\label{1 for constant coefficient}\\\\\n&&\\frac{d}{dt}\\langle (\\Delta z)^2\n\\rangle =2\\kappa_z \\langle (\\Delta z)\n\\rangle+2\\kappa_{zz}\n-2\\kappa_{tz}\\kappa_z.\n\\label{2 for constant coefficient}\n\\end{eqnarray}\nBy combining the latter two formulas\none can derive\n\\begin{eqnarray}\n\\kappa_{zz}=\\frac{1}{2}\n\\frac{d\\sigma^2}{dt}\n+\\kappa_{tz}\\kappa_z.\n\\label{center formula-3}\n\\end{eqnarray}\nWe can find\nthat DCDV\nis the special case of\nthe latter equation\nwhen $\\kappa_{tz}\\kappa_z=0$.\n\nBecause the coefficients\n$\\kappa_z$, $\\kappa_{zz}$\nand $\\kappa_{tz}$ are all\nconstants,\nEquation (\\ref{center formula-3})\ncan\nbe rewritten as\n\\begin{eqnarray}\n\\frac{d\\sigma^2}{dt}=2\\kappa_{zz}\n-2\\kappa_{tz}\\kappa_z=constant.\n\\label{sigma}\n\\end{eqnarray}\nBy setting the initial value\nof $\\sigma^2$ as zero, from\nthe latter equation we can\nobtain\n\\begin{eqnarray}\n\\sigma^2=\\alpha t.\n\\label{formula of diffusion regime}\n\\end{eqnarray}\nwith\n\\begin{eqnarray}\n\\alpha=2\\kappa_{tz}\\kappa_z\n+2\\kappa_{zz}=constant\n\\label{alpha}\n\\end{eqnarray}\nEquation\n(\\ref{formula of diffusion regime})\ndescribes diffusive regime.\nTherefore, for $D_{\\mu\\mu}\n=D_{\\mu\\mu}(\\mu)$\nthe constant coefficient linear\ndifferential equation of the\nisotropic distribution function\n(see Equation (\\ref{equation of F\nwith constant coefficient})) can\ndescribe the diffusion\nof the energetic particles.\nIf $\\kappa_z$ or $\\kappa_{tz}$\ncan be ignored,\nDCDV (Equation (\\ref{kzz-sigma}))\nis obtained.\n\n\\subsection{Case 2:\ndiffusion convection equation}\n\\label{For the case diffusion\nconvection equation}\nIf Equation\n(\\ref{equation of F with\nconstant coefficient})\nis simplified to the\ndiffusion convection equation\n\\begin{eqnarray}\n\\frac{\\partial{F}}{\\partial{t}}\n=-\\kappa_z\\frac{\\partial{F}}\n{\\partial{z}}+\\kappa_{zz}\n\\frac{\\partial^2{F}}{\\partial{z^2}},\n\\label{simplified equation of\nF with constant coefficient-1}\n\\end{eqnarray}\nEquation (\\ref{center formula-3})\nbecomes DCDV.\n\n\\subsection{Case 3:\ndiffusion equation}\n\\label{For the case diffusion equation}\nIf the convection term can be\nneglected, the diffusion convection\nEquation (\\ref{simplified equation of\n\tF with constant coefficient-1})\ncan be simplified\nto diffusion equation\n\\begin{eqnarray}\n\\frac{\\partial{F}}{\\partial{t}}\n=\\kappa_{zz}\\frac{\\partial^2{F}}\n{\\partial{z^2}}.\n\\end{eqnarray}\nTherefore, Equation\n(\\ref{1 for constant coefficient})\nbecomes\n\\begin{equation}\n\\frac{d}{dt}\\langle (\\Delta z) \\rangle=0.\n\\label{1-result for constant ceofficient}\n\\end{equation}\nConsequently,\nit is easy to see that\nthe parallel diffusion coefficient\ncan be shown as\n\\begin{eqnarray}\n\\kappa_{zz}=\\frac{1}{2}\\frac{d}{dt}\n\\langle (\\Delta z)^2 \\rangle .\n\\label{center formula corresponding\nto diffusion equation}\n\\end{eqnarray}\n\n\\subsection{Case 4: the constant\nbackground magnetic field}\n\\label{For the case: the constant\nbackground magnetic field}\nFor the constant background magnetic\nfield, the Fokker-Planck equation with\ntime-dependent pitch-angle diffusion\ncoefficient is\n\\begin{equation}\n\\frac{\\partial{f}}{\\partial{t}}\n+ v\\mu\n\\frac{\\partial{f}}\n{\\partial{z}}=\n\\frac{\\partial{}}{\\partial{\\mu}}\n\\left[D_{\\mu \\mu}(\\mu,t)\n\\frac{\\partial{f}}\n{\\partial{\\mu}}\\right].\n\\label{Fokker-Planck equation\nfor constant mean magnetic field}\n\\end{equation}\nBy using the similar method\nin Section 2,\nfrom the latter equation we can\nobtain the variable coefficient\ndifferential equation of isotropic\ndistribution function as\n\\begin{eqnarray}\n\\frac{\\partial{F}}{\\partial{t}}=\n&&\\left(\\kappa_{zz}(t)\n\\frac{\\partial^2{F}}{\\partial{z^2}}\n+\\kappa_{zzz}(t)\n\\frac{\\partial^3{F}}{\\partial{z^3}}\n+\\kappa_{zzzz}(t)\n\\frac{\\partial^4{F}}{\\partial{z^4}}\n+\\cdots\\right)\n+\\left( \\kappa_{tz}(t)\n\\frac{\\partial^2{F}}{\\partial{t}\\partial{z}}\n+ \\kappa_{ttz}(t)\n\\frac{\\partial^3{F}}\n{\\partial{t^2}\\partial{z}}\n+ \\kappa_{tttz}(t)\n\\frac{\\partial^4{F}}{\\partial{t^3}\n\\partial{z}}+\\cdots\\right)\\nonumber\\\\\n&&+\\left( \\kappa_{tzz}(t)\n\\frac{\\partial^3{F}}{\\partial{t}\n\\partial{z^2}}\n+ \\kappa_{ttzz}(t)\\frac{\\partial^4{F}}\n{\\partial{t^2}\\partial{z^2}}\n+ \\kappa_{tttzz}(t)\\frac{\\partial^5{F}}\n{\\partial{t^3}\\partial{z^2}}+\\cdots\\right)\n+\\cdots\\cdots.\n\\label{equation of F for\nconstant mean magnetic field}\n\\end{eqnarray}\n\nComparing\nwith Equation (\\ref{equation of F}),\nwe find that Equation\n(\\ref{equation of F for\nconstant mean magnetic field})\ndoes not contain the convection\nterm $\\kappa_z(t)\n\\partial{F}\/\\partial{z}$\nwhich is caused\nby the adiabatic\nfocusing effect.\nThe first order moment\nof parallel displacement\ncorresponding to Equation\n(\\ref{equation of F for\nconstant mean magnetic field})\nis given as,\n\\begin{eqnarray}\n&&\\frac{d}{dt}\\langle (\\Delta z)\n\\rangle=\\kappa_z(t)=0.\n\\label{1-result for constant\nmean magentic field}\n\\end{eqnarray}\nSimilarily, we can find the second\norder moment of the parallel\ndisplacement as\n\\begin{eqnarray}\n\\frac{d}{dt}\\langle (\\Delta z)^2\n\\rangle =2\\kappa_{zz}(t).\n\\label{2-result for constant\nmean magentic field}\n\\end{eqnarray}\nThis equation can be\nrewritten as\n\\begin{eqnarray}\n\\kappa_{zz}(t)=\\frac{1}{2}\n\\frac{d}{dt}\\langle (\\Delta z)^2\n\\rangle.\n\\label{kzz for constant mean\nmagentic field}\n\\end{eqnarray}\n\n\\section{SUMMARY AND CONCLUSION}\n\\label{SUMMARY AND CONCLUSION}\n\nThe\nparallel diffusion coefficient\ncan be defined as\none half of the temporal derivative\nof variance of the parallel displacement,\nDCDV (Equation (\\ref{kzz-sigma})).\nIt is shown that DCDV is only suitable to\nthe transport equation\nwithout the convection term or\nwithout the cross terms.\nIf the parallel diffusion coefficient\n$\\kappa_{zz}$ is a constant,\nthe formula $\\sigma^2=2\\kappa_{zz}t$ holds,\nwhich means the transport regime is diffusion.\nThe purpose of this work is to derive\nthe formula of the parallel\ndiffusion coefficient $\\kappa_{zz}$\n with the parallel displacement\n variance $\\sigma^2$\nfor the more general transport\nequation of the isotropic\ndistribution function.\nWe start from the Fokker-Planck equation\nwith time-dependent pitch-angle\ndiffusion coefficient $D_{\\mu\\mu}(\\mu,t)$,\nwhich is an extension of the\nFokker-Planck equation with\ntime-independent\npitch-angle diffusion coefficient.\nFrom this equation,\nby employing the\nperturbation method which is frequently\nused in relevant research\n\\citep[see, e.g.,][]{BeeckEA1986,\nBieberEA1990, SchlickeiserEA2008,\nSchlickeiserEA2010,\nLitvinenkoASchlickeiser2013,\nHeEA2014, WangEA2017b, wq2018},\nthe variable coefficient differential\nequation of the isotropic distribution\n function $F(z,t)$\nwith the time-dependent coefficients\nis derived.\nFrom the equation of $F(z,t)$ the modified\nformula describing the relation\nof the parallel diffusion coefficient\n$\\kappa_{zz}$ with the parallel\ndisplacement variance $\\sigma^2$\n (MDCDV)\nis found.\nMDCDV indicates that the parallel\ndiffusion is determined by not only\nthe parallel displacement variance\nbut also all\nthe coefficients of the cross terms\nwith the first order spatial derivative\nand convection term coefficient.\nBecause all coefficients are\ntime-dependent in MDCDV,\nthe well-known formula $\\sigma^2\\propto t$\ndoes not hold.\nTherefore, the transport of\nthe energetic charged particles described\nby MDCDV\nis not in diffusion regime.\n\nIn addition, we study\nMDCDV in some special cases.\nFor the isotropic distribution\nfunction equation with\nthe constant coefficient,\nMDCDV is simplified to\nEquation (\\ref{center formula-3}).\nThis formula demonstrates that\nthe parallel diffusion coefficient\nis determined by\nthe parallel displacement variance,\nthe convection term coefficient, and\nthe cross term coefficients with all the first\n order temporal and spatial\nderivatives.\nSince all the coefficients are constant,\n $\\sigma^2\\propto t$ can be found.\nTherefore, the constant\ncoefficient differential\nequation of $F(z,t)$ can\ndescribe transport in diffusion\nregime.\nFurthermore, we explore the\ntransport equation\nwith uniform background magnetic field.\nFor this case, the convection term does\n not occur in the equation of $F(z,t)$.\nAccordingly, MDCDV derived\nin this paper\nreduces to DCDV.\nFor the diffusion convection\nequation and\ndiffusion equation,\nthe more simplified equations of\n$F(z,t)$,\nMDCDV is also reduced to\nDCDV.\nThe telegraph equation is\nanother frequently used\ntransport equation,\nwhich is derived by using iteration\nmethod from the isotropic\ndistribution function equation.\nIn this article,\nwe also derive a new telegraph equation,\nand explore preliminarily the corresponding\nformula of $\\kappa_{zz}$\nwith variance $\\sigma^2$.\n\nInstead of adibatic focusing,\nsolar wind speed is another factor\ncausing the convection\nterm in the equation of the isotropic\ndistribution function.\nSo, when the phase space equation\nwith solar\nwind speed is projected to real space,\nsolar wind leads to a convection\nterm in the real space equation.\nWe\nwill explore this topic in the\nfuture.\nAs a subject of future work we will\nperform a\ndetailed investigation into\nthe new\ntelegraph equation.\n\n\\acknowledgments\n\nWe are partly supported by\ngrant NNSFC 41874206, NNSFC 41574172.\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nNotions of dimension abound in mathematics, and they give us quantitative measures of the sizes of various mathematical objects in a broad sense. In some instances, one wishes to know the exact dimension while in other instances, one just wishes to determine finiteness of the dimension. Finiteness of various dimensions has been considered in connection with central problems in the theory of $C^*$-algebras and in noncommutative geometry. For instance, finiteness of nuclear dimension (a noncommutative analog of covering dimension introduced in \\cite{WinZac}) plays a crucial role in the classification of $C^*$-algebras (e.g. \\cite{Ell,EGLN,EGLN17,Tik,TWW} and references therein), while finiteness of asymptotic dimension (a coarse geometric analog of covering dimension introduced in \\cite{Grom}) has featured in work on the Baum-Connes conjecture and the Novikov conjecture (e.g. \\cite{CG,CFY,Yu98}). In this paper, we will consider the notion of dynamical complexity and the implication of its finiteness on an $L^p$ assembly map, which is a model for the Baum-Connes assembly map with coefficients when $p=2$.\n\nDynamical complexity is a property of topological dynamical systems (and more generally, for \\'{e}tale groupoids) introduced by Guentner, Willett, and Yu in \\cite{GWY2}. Its definition was inspired by the notion of dynamic asymptotic dimension introduced in \\cite{GWY1} and by the notion of decomposition complexity introduced in \\cite{GTY}. Given an action of a countable discrete group $\\Gamma$ on a compact Hausdorff space $X$, the action is said to have finite dynamical complexity if the transformation groupoid $\\Gamma\\ltimes X$ is contained in the smallest class of open subgroupoids of $\\Gamma\\ltimes X$ that contains all relatively compact open subgroupoids and is closed under decomposability. Here, decomposability of an open subgroupoid $G$ over a collection $\\mathcal{C}$ of open subgroupoids roughly means that at any given scale, there is a cover of the unit space of $G$ by two open sets such that the subgroupoids associated to the two open sets at that scale are both in $\\mathcal{C}$. We refer the reader to Definition \\ref{FDCbasic} for the precise definition and to \\cite[Definition A.4]{GWY2} for a definition applicable to general \\'{e}tale groupoids.\n\n\nThe authors of \\cite{GWY2} considered a model for the Baum-Connes assembly map for an action based on equivariant versions of Yu's localization algebras (introduced in \\cite{Yu97}) and Roe algebras (introduced in \\cite{Roe88,Roe93}). In the appendix of \\cite{GWY2}, the authors showed that this model for the Baum-Connes assembly map is equivalent to the more traditional one stated in terms of Kasparov's $KK$-theory \\cite{BCH}. Their main result is the following:\n\n\\begin{thm}\\cite{GWY2}\nSuppose an action of a countable discrete group $\\Gamma$ on a compact Hausdorff space $X$ has finite dynamical complexity. Then $\\Gamma$ satisfies the Baum-Connes conjecture with coefficients in $C(X)$.\n\\end{thm}\n\n\nAlthough the aforementioned result follows from earlier work of Tu \\cite{Tu} on the Baum-Connes conjecture for amenable groupoids, the proof given in \\cite{GWY2} is completely different, and in some sense more direct and more elementary. In fact, their proof is inspired by Yu's proof of the coarse Baum-Connes conjecture for spaces with finite asymptotic dimension in \\cite{Yu98}. The main tool in both cases is a controlled Mayer-Vietoris sequence, which is part of a framework of quantitative (or controlled) $K$-theory for $C^*$-algebras developed by Yu together with Oyono-Oyono in \\cite{OY15,OY}. Roughly speaking, finite dimension\/complexity enables one to apply the Mayer-Vietoris argument a finite number of times to arrive at the quantitative $K$-theory of the algebra in question, and the (standard) $K$-theory of the algebra is obtained as a limit of the quantitative $K$-theory. \n\nIn our earlier work \\cite{Chung1}, we have extended the framework of quantitative $K$-theory to a larger class of Banach algebras, so that it can be applied to algebras of bounded linear operators on $L^p$ spaces in particular. Our goal in this paper is to consider the $L^p$ analog of the assembly map in \\cite{GWY2}, and use our extended framework of quantitative $K$-theory to show that this assembly map is an isomorphism under the assumption of finite dynamical complexity. In fact, one sees that the techniques and proofs in \\cite{GWY2} carry over to our setting with minor adjustments, the main differences being the exposition of the homotopy invariance argument in Section \\ref{sec:Base}, and more care in the use of quantitative $K$-theory (due to an additional norm control parameter) in the Mayer-Vietoris argument used in the proof of the main theorem in Section \\ref{sec:Ind}. \n\nIn order to explain our setup, let us briefly recall the Baum-Connes conjecture with coefficients.\nGiven a (separable) $C^*$-algebra $A$ and an action of a countable discrete group $\\Gamma$ on $A$ by $*$-automorphisms, one may form the reduced crossed product $C^*$-algebra $A\\rtimes_\\lambda\\Gamma$. The Baum-Connes conjecture with coefficients \\cite{BCH} posits that a certain homomorphism (or assembly map) \\[\\mu:K_*^\\Gamma(\\underline{E}\\Gamma;A)\\rightarrow K_*(A\\rtimes_\\lambda\\Gamma)\\] is an isomorphism, where the left-hand side is the equivariant $K$-homology with coefficients in $A$ of the classifying space $\\underline{E}\\Gamma$ for proper $\\Gamma$-actions, and the right-hand side is the $K$-theory of the reduced crossed product $C^*$-algebra. Consider a particular model for $\\underline{E}\\Gamma$, namely $\\bigcup_{s\\geq 0}P_s(\\Gamma)$ equipped with the $\\ell^1$ metric (cf. \\cite[Section 2]{BCH}), where $P_s(\\Gamma)$ is the Rips complex of $\\Gamma$ at scale $s$, i.e., it is the simplicial complex with vertex set $\\Gamma$, and where a finite subset $E\\subset\\Gamma$ spans a simplex if and only if $d(g,h)\\leq s$ for all $g,h\\in E$. Here we assume that $\\Gamma$ is equipped with a proper length function and $d$ is the associated metric. One may then reformulate the Baum-Connes assembly map as \\[\\lim_{s\\rightarrow\\infty}K_*(C_L^*(P_s(\\Gamma);A))\\stackrel{\\epsilon_0}{\\rightarrow} \\lim_{s\\rightarrow\\infty}K_*(C^*(P_s(\\Gamma);A))\\cong K_*(A\\rtimes_\\lambda\\Gamma),\\]\nwhere $C^*(P_s(\\Gamma);A)$ is the equivariant Roe algebra with coefficients in $A$, $C_L^*(P_s(\\Gamma);A)$ is Yu's localization algebra with coefficients in $A$, and $\\epsilon_0$ is induced by the evaluation-at-zero map. The fact that $K$-homology can be identified with the $K$-theory of the localization algebra was shown for finite-dimensional simplicial complexes in \\cite{Yu97}, and for general locally compact metric spaces in \\cite{QR}. The fact that the equivariant Roe algebra with coefficients is stably isomorphic to the reduced crossed product underlies the coarse-geometric approach to the Baum-Connes conjecture with coefficients. We refer the reader to \\cite[Appendix B]{GWY2} for a detailed comparison of the abovementioned assembly maps, and also \\cite{Roe02} for a comparison of assembly maps without coefficients.\n\nNow let $A$ be a norm-closed subalgebra of $B(L^p(Z,\\mu))$ for some measure space $(Z,\\mu)$ and $p\\in[1,\\infty)$. We refer to such algebras as $L^p$ operator algebras. Also suppose that a countable discrete group $\\Gamma$ acts on $A$ by isometric automorphisms. By mimicking the $C^\\ast$-algebraic definitions, we can define $L^p$ analogs $B^p_L(P_s(\\Gamma);A)$, $B^p(P_s(\\Gamma);A)$, $A\\rtimes_{\\lambda,p}\\Gamma$ of the localization algebra, the Roe algebra, and the reduced crossed product respectively. We can then consider the following $L^p$ assembly map for the $\\Gamma$-action:\n\\[\\lim_{s\\rightarrow\\infty}K_*(B^p_L(P_s(\\Gamma);A))\\stackrel{\\epsilon_0}{\\rightarrow} \\lim_{s\\rightarrow\\infty}K_*(B^p(P_s(\\Gamma);A))\\cong K_*(A\\rtimes_{\\lambda,p}\\Gamma).\\]\nIn this paper, $A$ will be $C(X)$, where $X$ is a compact Hausdorff space, and\nour main result may be stated as follows:\n\\begin{thm} (cf. Theorem \\ref{mainthm}) \\label{LpBC}\nSuppose that an action of a countable discrete group $\\Gamma$ on a compact Hausdorff space $X$ has finite dynamical complexity. Then the $L^p$ assembly map for the action is an isomorphism for $p\\in[1,\\infty)$.\n\\end{thm}\n\nA result like this can be seen as an indication of the computability of the $K$-theory of the $L^p$ reduced crossed product on the right-hand side of the assembly map since the $K$-theory of localization-type algebras like the one on the left-hand side of the assembly map has Mayer-Vietoris sequences associated to decompositions of the simplicial complex and other properties of a generalized homology theory (cf. \\cite{Yu97,CW}). Moreover, it enables one to transfer questions about whether the $K$-theory of certain $L^p$ reduced crossed products is independent of $p$ (cf. \\cite[Problem 11.2]{PhilOpen} and \\cite{LY}) over to the left-hand side of the assembly map. On a different note, our method of proof indicates that arguments involving quantitative (or controlled) $K$-theory in the $C^\\ast$-algebraic setting may be adapted without much difficulty to the $L^p$ setting, while there has been substantial difficulty in carrying over other approaches to the Baum-Connes conjecture, such as the Dirac-dual Dirac method.\n\n\n\nIn Section 2, we define $L^p$ Roe algebras and localization algebras associated to an action of a countable discrete group on a compact Hausdorff space, and we define the $L^p$ assembly map in terms of the $K$-theory of these algebras. In the case $p=2$, these are exactly the algebras and map considered in \\cite{GWY2}. In Section 3, we associate subalgebras of these algebras to subgroupoids of the transformation groupoid given by the action, and recall the notion of dynamical complexity. In Section 4, we recall some definitions and facts from the framework of quantitative $K$-theory that we developed in \\cite{Chung1}. Finally, in Section 5, we prove our main result via a homotopy invariance argument and a Mayer-Vietoris argument.\n\n\n\n\\section{An $L^p$ assembly map} \\label{Sect:assembly}\n\nThroughout this section, $\\Gamma$ will be a countable discrete group acting on a compact Hausdorff space $X$ by homeomorphisms. The action will be denoted by $\\Gamma\\curvearrowright X$. We also assume that $\\Gamma$ is equipped with a proper length function $l:\\Gamma\\rightarrow\\mathbb{N}$ and the associated right invariant metric. We will define an assembly map in terms of $L^p$ versions of localization algebras and Roe algebras, where $p\\in[1,\\infty)$. When $p=2$, we recover (a model for) the Baum-Connes assembly map for $\\Gamma$ with coefficients in $C(X)$ considered in \\cite{GWY2}.\n\n\n\\begin{defn}\nLet $s\\geq 0$. The Rips complex of $\\Gamma$ at scale $s$, denoted $P_s(\\Gamma)$, is the simplicial complex with vertex set $\\Gamma$, and where a finite subset $E\\subset\\Gamma$ spans a simplex if and only if $d(g,h)\\leq s$ for all $g,h\\in E$.\n\nPoints in $P_s(\\Gamma)$ can be written as formal linear combinations $\\sum_{g\\in\\Gamma}t_g g$, where $t_g\\in[0,1]$ for each $g$ and $\\sum_{g\\in\\Gamma}t_g=1$. We equip $P_s(\\Gamma)$ with the $\\ell^1$ metric, i.e., $d(\\sum_{g\\in\\Gamma}t_g g,\\sum_{g\\in\\Gamma}s_g g)=\\sum_{g\\in\\Gamma}|t_g-s_g|$.\n\nThe barycentric coordinates on $P_s(\\Gamma)$ are the continuous functions \\[t_g:P_s(\\Gamma)\\rightarrow[0,1]\\] uniquely determined by the condition $z=\\sum_{g\\in\\Gamma}t_g(z)g$ for all $z\\in P_s(\\Gamma)$.\n\\end{defn}\n\nProperness of the length function on $\\Gamma$ implies that $P_s(\\Gamma)$ is finite dimensional and locally compact. Also, the right translation action of $\\Gamma$ on itself extends to a right action of $\\Gamma$ on $P_s(\\Gamma)$ by isometric simplicial automorphisms.\n\nIn the usual setting of the Baum-Connes conjecture (e.g. in \\cite{GWY2}), one considers Hilbert spaces and $C^*$-algebras encoding the large scale geometry of $\\Gamma$ and the topology of $P_s(\\Gamma)$. We will replace these Hilbert spaces by $L^p$ spaces, thereby obtaining $L^p$ operator algebras instead of $C^*$-algebras.\n\n\nFirst, we recall some facts about $L^p$ tensor products. Details can be found in \\cite[Chapter 7]{DF}.\n\nFor $p\\in[1,\\infty)$, there is a tensor product of $L^p$ spaces such that we have a canonical isometric isomorphism $L^p(X,\\mu)\\otimes L^p(Y,\\nu)\\cong L^p(X\\times Y,\\mu\\times\\nu)$, which identifies, for every $\\xi\\in L^p(X,\\mu)$ and $\\eta\\in L^p(Y,\\nu)$, the element $\\xi\\otimes\\eta$ with the function $(x,y)\\mapsto\\xi(x)\\eta(y)$ on $X\\times Y$. Moreover, this tensor product has the following properties:\n\\begin{enumerate}\n\\item Under the identification above, the linear span of all $\\xi\\otimes\\eta$ is dense in $L^p(X\\times Y,\\mu\\times\\nu)$.\n\\item $||\\xi\\otimes\\eta||_p=||\\xi||_p||\\eta||_p$ for all $\\xi\\in L^p(X,\\mu)$ and $\\eta\\in L^p(Y,\\nu)$.\n\\item The tensor product is commutative and associative.\n\\item If $a\\in B(L^p(X_1,\\mu_1),L^p(X_2,\\mu_2))$ and $b\\in B(L^p(Y_1,\\nu_1),L^p(Y_2,\\nu_2))$, then there exists a unique \\[c\\in B(L^p(X_1\\times Y_1,\\mu_1\\times\\nu_1),L^p(X_2\\times Y_2,\\mu_2\\times\\nu_2))\\] such that under the identification above, $c(\\xi\\otimes\\eta)=a(\\xi)\\otimes b(\\eta)$ for all $\\xi\\in L^p(X_1,\\mu_1)$ and $\\eta\\in L^p(Y_1,\\nu_1)$. We will denote this operator by $a\\otimes b$. Moreover, $||a\\otimes b||=||a|| ||b||$.\n\\item The tensor product of operators is associative, bilinear, and satisfies $(a_1\\otimes b_1)(a_2\\otimes b_2)=a_1a_2\\otimes b_1b_2$.\n\\end{enumerate}\nIf $A\\subseteq B(L^p(X,\\mu))$ and $B\\subseteq B(L^p(Y,\\nu))$ are norm-closed subalgebras, we then define $A\\otimes B\\subseteq B(L^p(X\\times Y,\\mu\\times\\nu))$ to be the closed linear span of all elements of the form $a\\otimes b$ with $a\\in A$ and $b\\in B$.\n\n\n\\begin{defn}\nFor $s\\geq 0$, define \\[Z_s=\\biggl\\{\\sum_{g\\in\\Gamma}t_g g\\in P_s(\\Gamma):t_g\\in\\mathbb{Q}\\;\\text{for all}\\;g\\in\\Gamma\\biggr\\}.\\]\nNote that $Z_s$ is a $\\Gamma$-invariant, countable, dense subset of $P_s(\\Gamma)$.\n\nDefine \\[E_s=\\ell^p(Z_s)\\otimes\\ell^p(X)\\otimes\\ell^p\\otimes\\ell^p(\\Gamma)\\cong \\ell^p(Z_s\\times X,\\ell^p\\otimes\\ell^p(\\Gamma)),\\] and equip $E_s$ with the isometric $\\Gamma$-action given by \\[u_g\\cdot(\\delta_z\\otimes\\delta_x\\otimes\\eta\\otimes\\delta_h)=\\delta_{zg^{-1}}\\otimes\\delta_{gx}\\otimes\\eta\\otimes\\delta_{gh}\\] for $z\\in Z_s$, $x\\in X$, $\\eta\\in\\ell^p$, and $g,h\\in\\Gamma$. \n\\end{defn}\n\nNote that if $s_0\\leq s$, then $P_{s_0}(\\Gamma)$ identifies equivariantly and isometrically with a subcomplex of $P_s(\\Gamma)$, and $Z_{s_0}\\subset Z_s$. Hence we have a canonical equivariant isometric inclusion $E_{s_0}\\subset E_s$.\n\nWe will write $\\mathcal{K}_\\Gamma$ for the algebra of compact operators on $\\ell^p\\otimes\\ell^p(\\Gamma)\\cong\\ell^p(\\mathbb{N}\\times\\Gamma)$ equipped with the $\\Gamma$-action induced by the tensor product of the trivial action on $\\ell^p$ and the left regular representation on $\\ell^p(\\Gamma)$. We also equip the algebra $C(X)\\otimes \\mathcal{K}_\\Gamma$ with the diagonal action of $\\Gamma$. Note that the natural faithful representation of $C(X)\\otimes \\mathcal{K}_\\Gamma$ on $\\ell^p(X)\\otimes\\ell^p\\otimes\\ell^p(\\Gamma)$ is covariant for the representation defined by tensoring the natural action on $\\ell^p(X)$, the trivial representation on $\\ell^p$, and the regular representation on $\\ell^p(\\Gamma)$. \n\nNow we can define the $L^p$ operator algebras that will feature in our assembly map.\n\n\\begin{defn}\nLet $T$ be a bounded linear operator on $E_s$, which we may regard as a ($Z_s\\times Z_s$)-indexed matrix $T=(T_{y,z})$ with \\[T_{y,z}\\in B(\\ell^p(X)\\otimes\\ell^p\\otimes\\ell^p(\\Gamma))\\] for each $y,z\\in Z_s$.\n\\begin{enumerate}\n\\item $T$ is $\\Gamma$-invariant if $u_g T u_g^{-1}=T$ for all $g\\in\\Gamma$.\n\\item The Rips-propagation of $T$ is $\\sup\\{d_{P_s(\\Gamma)}(y,z):T_{y,z}\\neq 0\\}$.\n\\item The $\\Gamma$-propagation of $T$, denoted by $\\mathrm{prop}_\\Gamma(T)$, is \\[\\sup\\{d_\\Gamma(g,h):T_{y,z}\\neq 0\\;\\text{for some}\\; y,z\\in Z_s\\;\\text{with}\\; t_g(y)\\neq 0\\;\\text{and}\\;t_h(z)\\neq 0\\}.\\] \n\\item $T$ is $X$-locally compact if $T_{y,z}\\in C(X)\\otimes \\mathcal{K}_\\Gamma$ for all $y,z\\in Z_s$, and if for any compact subset $F\\subset P_s(\\Gamma)$, the set \\[\\{(y,z)\\in (F\\times F)\\cap(Z_s\\times Z_s):T_{y,z}\\neq 0\\}\\] is finite.\n\\end{enumerate}\n\\end{defn} \n\n\n\\begin{defn}\nLet $\\mathbb{C}[\\Gamma\\curvearrowright X;s]$ denote the algebra of all $\\Gamma$-invariant, $X$-locally compact operators on $E_s$ with finite $\\Gamma$-propagation.\n\nLet $B^p(\\Gamma\\curvearrowright X;s)$ denote the closure of $\\mathbb{C}[\\Gamma\\curvearrowright X;s]$ with respect to the operator norm on $B(E_s)$. We will call $B^p(\\Gamma\\curvearrowright X;s)$ the (equivariant) $L^p$ Roe algebra of $\\Gamma\\curvearrowright X$ at scale $s$.\n\\end{defn}\n\n\nWe will regard the algebras above as concretely represented on $E_s$, and we will often think of elements of $B^p(\\Gamma\\curvearrowright X;s)$ as matrices $(T_{y,z})_{y,z\\in Z_s}$ with entries being continuous equivariant functions $T_{y,z}:X\\rightarrow \\mathcal{K}_\\Gamma$ (with additional properties).\n\n\n\\begin{defn}\nLet $\\mathbb{C}_L[\\Gamma\\curvearrowright X;s]$ denote the algebra of all bounded, uniformly continuous functions $a:[0,\\infty)\\rightarrow\\mathbb{C}[\\Gamma\\curvearrowright X;s]$ such that the $\\Gamma$-propagation of $a(t)$ is uniformly finite as $t$ varies, and such that the Rips-propagation of $a(t)$ tends to zero as $t\\rightarrow\\infty$. \n\nLet $B^p_L(\\Gamma\\curvearrowright X;s)$ denote the completion of $\\mathbb{C}_L[\\Gamma\\curvearrowright X;s]$ with respect to the norm \\[||a||:=\\sup_{t\\in[0,\\infty)} ||a(t)||_{B^p(\\Gamma\\curvearrowright X;s)}.\\]\nWe will call $B^p_L(\\Gamma\\curvearrowright X;s)$ the $L^p$ localization algebra of $\\Gamma\\curvearrowright X$ at scale $s$.\n\\end{defn}\n\nWe will regard $B^p_L(\\Gamma\\curvearrowright X;s)$ as concretely represented on $L^p[0,\\infty)\\otimes E_s$, and elements of $B^p_L(\\Gamma\\curvearrowright X;s)$ can be regarded as bounded, uniformly continuous functions $a:[0,\\infty)\\rightarrow B^p(\\Gamma\\curvearrowright X;s)$ (with additional properties).\n\nNow consider the evaluation-at-zero homomorphism \\[\\epsilon_0:B^p_L(\\Gamma\\curvearrowright X;s)\\rightarrow B^p(\\Gamma\\curvearrowright X;s),\\] which induces a homomorphism on $K$-theory \\[\\epsilon_0:K_*(B^p_L(\\Gamma\\curvearrowright X;s))\\rightarrow K_*(B^p(\\Gamma\\curvearrowright X;s)).\\]\n\nIf $s_0\\leq s$, then the equivariant isometric inclusion $E_{s_0}\\subset E_s$ allows us to regard $\\mathbb{C}[\\Gamma\\curvearrowright X;s_0]$ as a subalgebra of $\\mathbb{C}[\\Gamma\\curvearrowright X;s]$. We then regard $B^p(\\Gamma\\curvearrowright X;s_0)$ (resp. $B^p_L(\\Gamma\\curvearrowright X;s_0)$) as a subalgebra of $B^p(\\Gamma\\curvearrowright X;s)$ (resp. $B^p_L(\\Gamma\\curvearrowright X;s)$). Thus there are directed systems of inclusions of $L^p$ operator algebras $(B^p(\\Gamma\\curvearrowright X;s))_{s\\geq 0}$ and $(B^p_L(\\Gamma\\curvearrowright X;s))_{s\\geq 0}$, and the evaluation-at-zero maps above are compatible with these inclusions.\n\n\\begin{defn} \\label{assembly}\nThe $L^p$ assembly map for the action $\\Gamma\\curvearrowright X$ is the direct limit\n\\[ \\epsilon_0:\\lim_{s\\rightarrow\\infty}K_*(B^p_L(\\Gamma\\curvearrowright X;s))\\rightarrow\\lim_{s\\rightarrow\\infty}K_*(B^p(\\Gamma\\curvearrowright X;s)). \\]\n\\end{defn}\nWhen $p=2$, this is the model for the Baum-Connes assembly map considered in \\cite{GWY2}.\n\nRegarding $A=C(X)$ as an $L^p$ operator algebra (acting by multiplication on $L^p(X,\\mu)$ for some measure $\\mu$) with an isometric $\\Gamma$-action, we can define an $L^p$ reduced crossed product as follows:\n\nConsider $C_c(\\Gamma,A)$, the set of finite sums of the form $\\sum_{g\\in\\Gamma}a_g g$ with $a_g\\in A$ and with the product given by \\[ \\biggl(\\sum_{g\\in\\Gamma}a_g g\\biggr)\\biggl(\\sum_{h\\in\\Gamma}b_h h\\biggr)=\\sum_{g,h\\in\\Gamma}a_g\\alpha_g(b_h)gh, \\] where $\\alpha$ denotes the $\\Gamma$-action on $A$. There is a natural faithful representation of $C_c(\\Gamma,A)$ on $\\ell^p(\\Gamma,L^p(X,\\mu))$ given by\n\\begin{align*}\n(a\\xi)(h) &= \\alpha_{h^{-1}}(a)\\xi(h), \\\\\n(g\\xi)(h) &= \\xi(g^{-1}h)\n\\end{align*}\nfor $a\\in A$, $g,h\\in\\Gamma$, and $\\xi\\in\\ell^p(\\Gamma,L^p(X,\\mu))$. We then define the $L^p$ reduced crossed product $A\\rtimes_{\\lambda,p}\\Gamma$ to be the operator norm closure of $C_c(\\Gamma,A)$ in $B(\\ell^p(\\Gamma,L^p(X,\\mu)))$.\n\nAdapting arguments in the $C^\\ast$-algebraic setting (cf. \\cite{GWY2} or \\cite[Lemma 3.7]{WY}), one can show that the equivariant $L^p$ Roe algebra $B^p(\\Gamma\\curvearrowright X;s)$ is stably isomorphic to the crossed product $C(X)\\rtimes_{\\lambda,p}\\Gamma$. In particular, the right-hand side of the $L^p$ assembly map defined above is the same as the $K$-theory of $C(X)\\rtimes_{\\lambda,p}\\Gamma$.\n\nFor most of the rest of this paper, we will work with the kernel of this $L^p$ assembly map.\n\n\\begin{defn}\nLet $B^p_{L,0}(\\Gamma\\curvearrowright X;s)$ be the subalgebra of $B^p_L(\\Gamma\\curvearrowright X;s)$ consisting of functions $a$ such that $a(0)=0$. We will call $B^p_{L,0}(\\Gamma\\curvearrowright X;s)$ the $L^p$ obstruction algebra of $\\Gamma\\curvearrowright X$ at scale $s$.\n\\end{defn}\n\n\\begin{lem}\\label{Lem:obsvan}\nThe $L^p$ assembly map for $\\Gamma\\curvearrowright X$ is an isomorphism if and only if \\[\\lim_{s\\rightarrow\\infty}K_*(B^p_{L,0}(\\Gamma\\curvearrowright X;s))=0.\\]\n\\end{lem}\n\n\\begin{proof}\nJust as in the $C^\\ast$-algebraic setting, we have a short exact sequence \\[0\\rightarrow B^p_{L,0}(\\Gamma\\curvearrowright X;s)\\rightarrow B^p_L(\\Gamma\\curvearrowright X;s)\\rightarrow B^p(\\Gamma\\curvearrowright X;s)\\rightarrow 0,\\] which induces the usual six-term exact sequence in $K$-theory. The lemma then follows from continuity of $K$-theory under direct limits, and the preservation of exact sequences under direct limits of abelian groups.\n\\end{proof}\n\nOur goal in this paper will be to show that if $\\Gamma\\curvearrowright X$ has finite dynamical complexity, then $\\lim_{s\\rightarrow\\infty}K_*(B^p_{L,0}(\\Gamma\\curvearrowright X;s))=0$, and thus the $L^p$ assembly map for $\\Gamma\\curvearrowright X$ in Definition \\ref{assembly} is an isomorphism.\n\n\n\n\n\n\n\\section{Groupoids and dynamical complexity}\n\nIn this section, we consider the transformation groupoid associated to a group action, its subgroupoids, and $L^p$ operator algebras associated with them. We also recall the definition of dynamical complexity from \\cite{GWY2}.\n\n\\begin{defn}\nThe transformation groupoid $\\Gamma\\ltimes X$ associated to $\\Gamma\\curvearrowright X$ is $\\{(gx,g,x):g\\in\\Gamma,x\\in X\\}$ topologized such that the projection $\\Gamma\\ltimes X\\rightarrow\\Gamma\\times X$ onto the second and third factors is a homeomorphism, and equipped with the following additional structure:\n\\begin{enumerate}\n\\item A pair $((hy,h,y),(gx,g,x))$ of elements in $\\Gamma\\ltimes X$ is said to be composable if $y=gx$. In this case, their product is defined by \\[(hgx,h,gx)(gx,g,x)=(hgx,hg,x).\\]\n\\item The inverse of an element $(gx,g,x)\\in\\Gamma\\ltimes X$ is \\[(gx,g,x)^{-1}=(x,g^{-1},gx).\\]\n\\item The units of $\\Gamma\\ltimes X$ are the elements of the clopen subspace \\[G^{(0)}=\\{(x,e,x):x\\in X\\},\\] where $e$ is the identity in $\\Gamma$. We call $G^{(0)}$ the unit space of $\\Gamma\\ltimes X$.\n\\end{enumerate}\n\\end{defn}\n\n\\begin{defn}\nLet $s\\geq 0$, and let $P_s(\\Gamma)$ be the Rips complex of $\\Gamma$ at scale $s$. The support of $z=\\sum_{g\\in\\Gamma}t_g(z) g\\in P_s(\\Gamma)$ is the finite set \\[\\mathrm{supp}(z)=\\{g\\in\\Gamma:t_g(z)\\neq 0\\}.\\]\nThe support of $T=(T_{y,z})_{y,z\\in Z_s}\\in B^p(\\Gamma\\curvearrowright X;s)$ is \n\\[ \\mathrm{supp}(T)=\\left\\{(gx,gh^{-1},hx)\\in\\Gamma\\ltimes X:\\;\\parbox[c][4em][c]{0.5\\textwidth}{there exist $y,z\\in Z_s$ with $T_{y,z}(x)\\neq 0$, $g\\in \\mathrm{supp}(y)$, and $h\\in \\mathrm{supp}(z)$}\\right\\}. \\]\n\\end{defn}\n\nWith this definition, one sees that \\[\\mathrm{prop}_\\Gamma(T)=\\sup\\{l(gh^{-1}):(gx,gh^{-1},hx)\\in \\mathrm{supp}(T)\\;\\text{for some}\\;x\\in X\\}.\\]\n\nGiven two subsets $A,B\\subset\\Gamma\\ltimes X$, we write $AB$ for \\[\\{ab:a\\in A,b\\in B,(a,b)\\;\\text{is composable}\\}.\\] With this notation, the following lemma says that supports of operators in $B^p(\\Gamma\\curvearrowright X;s)$ behave as expected under composition of operators.\n\n\\begin{lem} \\label{supplem}\nLet $S,T\\in B^p(\\Gamma\\curvearrowright X;s)$. Then $\\mathrm{supp}(ST)\\subseteq \\mathrm{supp}(S)\\mathrm{supp}(T)$.\n\\end{lem}\n\n\\begin{proof}\nSuppose that $(gx,gh^{-1},hx)\\in \\mathrm{supp}(ST)$. Then there are $y,z\\in Z_s$ such that $(ST)_{y,z}(x)\\neq 0$, $g\\in \\mathrm{supp}(y)$, and $h\\in \\mathrm{supp}(z)$. Thus there is $w\\in Z_s$ such that $S_{y,w}(x)\\neq 0$ and $T_{w,z}(x)\\neq 0$. If $k\\in \\mathrm{supp}(w)$, then $(gx,gk^{-1},kx)\\in \\mathrm{supp}(S)$ and $(kx,kh^{-1},hx)\\in \\mathrm{supp}(T)$, so $(gx,gh^{-1},hx)=(gx,gk^{-1},kx)(kx,kh^{-1},hx)\\in \\mathrm{supp}(S)\\mathrm{supp}(T)$.\n\\end{proof}\n\n\\begin{defn}\nLet $\\Gamma\\ltimes X$ be the transformation groupoid associated to $\\Gamma\\curvearrowright X$. A subgroupoid of $\\Gamma\\ltimes X$ is a subset $G\\subset \\Gamma\\ltimes X$ that is closed under composition, taking inverses, and units, i.e.,\n\\begin{enumerate}\n\\item If $(hgx,h,gx)$ and $(gx,g,x)$ are in $G$, then so is $(hgx,hg,x)$.\n\\item If $(gx,g,x)\\in G$, then $(gx,g,x)^{-1}\\in G$.\n\\item If $(gx,g,x)\\in G$, then $(x,e,x)\\in G$ and $(gx,e,gx)\\in G$, where $e$ is the identity in $\\Gamma$.\n\\end{enumerate}\nA subgroupoid is equipped with the subspace topology from $\\Gamma\\ltimes X$.\n\\end{defn}\n\nSubgroupoids of $\\Gamma\\ltimes X$ give rise to subalgebras of the Roe algebra, localization algebra, and obstruction algebra that we defined in the previous section.\n\n\\begin{lem}\nLet $G$ be an open subgroupoid of $\\Gamma\\ltimes X$. Define $\\mathbb{C}[G;s]$ to be the subspace of $\\mathbb{C}[\\Gamma\\curvearrowright X;s]$ consisting of all operators $T$ with support contained in a compact subset of $G$. Then $\\mathbb{C}[G;s]$ is a subalgebra of $\\mathbb{C}[\\Gamma\\curvearrowright X;s]$.\n\\end{lem}\n\n\\begin{proof}\nGiven Lemma \\ref{supplem}, it suffices to show that if $A$ and $B$ are two relatively compact subsets of $G$, then so is $AB$. To see this, first suppose that $A$ and $B$ are compact. Then any net in $AB$ has a convergent subnet since nets in $A$ and nets in $B$ have this property, and so $AB$ is compact. Now if $A$ and $B$ are relatively compact, then since $AB\\subset \\bar{A}\\bar{B}$ and $\\bar{A}\\bar{B}$ is compact, it follows that $AB$ is relatively compact.\n\\end{proof}\n\n\\begin{defn}\nLet $G$ be an open subgroupoid of $\\Gamma\\ltimes X$. Let $\\mathbb{C}_L[G;s]$ denote the subalgebra of $\\mathbb{C}_L[\\Gamma\\curvearrowright X;s]$ consisting of functions $a$ such that $\\bigcup_{t\\in[0,\\infty)} \\mathrm{supp}(a(t))$ has compact closure in $G$. \n\nLet $\\mathbb{C}_{L,0}[G;s]$ denote the ideal of $\\mathbb{C}_L[G;s]$ consisting of functions $a$ such that $a(0)=0$.\n\nLet $B^p(G;s)$, $B^p_L(G;s)$, and $B^p_{L,0}(G;s)$ denote the respective closures of $\\mathbb{C}[G;s]$, $\\mathbb{C}_L[G;s]$, and $\\mathbb{C}_{L,0}[G;s]$ in $B^p(\\Gamma\\curvearrowright X;s)$, $B^p_L(\\Gamma\\curvearrowright X;s)$, and $B^p_{L,0}(\\Gamma\\curvearrowright X;s)$.\n\\end{defn}\n\nNote that when $G=\\Gamma\\ltimes X$, we have $B^p(G;s)=B^p(\\Gamma\\curvearrowright X;s)$, and similarly for the localization algebra and obstruction algebra.\n\nSince we will be working mostly with the obstruction algebras, we introduce the following shorthand notation for these algebras. We also need to construct filtrations on these algebras in order to apply quantitative $K$-theory in the proof of the main theorem.\n\n\\begin{defn}\nLet $G$ be an open subgroupoid of $\\Gamma\\ltimes X$, and let $s\\geq 0$. Set $A^s(G)$ to be $B^p_{L,0}(G;s)$. For $r\\geq 0$, define \\[A^s(G)_r=\\{a\\in\\mathbb{C}_{L,0}[G;s]:\\mathrm{prop}_\\Gamma(a(t))\\leq r\\;\\text{for all}\\;t\\},\\] which is a linear subspace of $A^s(G)$.\n\nWhen $G=\\Gamma\\ltimes X$, we will simply write $A^s$ and $A^s_r$.\n\\end{defn}\n\n\\begin{lem} \\label{filtlem}\nLet $G$ be an open subgroupoid of $\\Gamma\\ltimes X$, and let $s\\geq 0$. Then the family $(A^s(G))_{r\\geq 0}$ of subspaces of $A^s(G)$ satisfies:\n\\begin{enumerate}\n\\item if $r_1\\leq r_2$, then $A^s(G)_{r_1}\\subset A^s(G)_{r_2}$;\n\\item $A^s(G)_{r_1}A^s(G)_{r_2}\\subset A^s(G)_{r_1+r_2}$ for all $r_1,r_2\\geq 0$;\n\\item $\\bigcup_{r\\geq 0}A^s(G)_r$ is dense in $A^s(G)$.\n\\end{enumerate}\n\\end{lem}\n\n\\begin{proof}\nNote that $a\\in A^s(G)_r$ if and only if \n\\begin{itemize}\n\\item $a\\in\\mathbb{C}_{L,0}[G;s]$, and \n\\item $l(g)\\leq r$ whenever $(gx,g,x)\\in \\mathrm{supp}(a(t))$ for some $t\\geq 0$.\n\\end{itemize}\nProperties (i) and (iii) follow immediately. \n\nFor (ii), if $a\\in A^s(G)_{r_1}$, $b\\in A^s(G)_{r_2}$, and $(gx,g,x)\\in \\mathrm{supp}(a(t)b(t))$ for some $t$, then by Lemma \\ref{supplem}, $(gx,g,x)=(gx,gh^{-1},hx)(hx,h,x)$ for some $(gx,gh^{-1},hx)\\in \\mathrm{supp}(a(t))$ and $(hx,h,x)\\in \\mathrm{supp}(b(t))$. Thus $l(g)\\leq l(gh^{-1})+l(h)\\leq r_1+r_2$ so $ab\\in A^s(G)_{r_1+r_2}$.\n\\end{proof}\n\n\n\nNote that if $S$ is an open subset of $\\Gamma\\ltimes X$, then $S$ generates an open subgroupoid of $\\Gamma\\ltimes X$ (cf. \\cite[ Lemma 5.2]{GWY1}).\n\n\\begin{defn}\nLet $G$ be an open subgroupoid of $\\Gamma\\ltimes X$, let $H$ be an open subgroupoid of $G$, and let $r\\geq 0$. The expansion of $H$ by $r$ relative to $G$, denoted by $H^{+r}$, is the open subgroupoid of $\\Gamma\\ltimes X$ generated by \\[H\\cup\\{(gx,g,x)\\in G:x\\in H^{(0)},l(g)\\leq r\\}.\\]\n\\end{defn}\n\nNote that $H^{+r}$ depends on $G$ although we do not indicate this in the notation.\n\n\\begin{lem} \\label{pdtlem}\nLet $G$ be an open subgroupoid of $\\Gamma\\ltimes X$, let $H$ be an open subgroupoid of $G$, and let $r,s\\geq 0$. Then \\[A^s(H)\\cdot A^s_r(G)\\cup A^s_r(G)\\cdot A^s(H)\\subseteq A^s(H^{+r}).\\]\n\\end{lem}\n\n\\begin{proof}\nFollows from Lemma \\ref{supplem}.\n\\end{proof}\n\n\\begin{lem}\nLet $G$ be an open subgroupoid of $\\Gamma\\ltimes X$, let $H$ be an open subgroupoid of $G$, and let $r_1,r_2\\geq 0$. Then $(H^{+r_1})^{+r_2}\\subseteq H^{+(r_1+r_2)}$.\n\\end{lem}\n\n\\begin{proof}\nIt suffices to show that \\[\\{(gx,g,x)\\in G:x\\in (H^{+r_1})^{(0)},l(g)\\leq r_2\\}\\subseteq H^{+(r_1+r_2)}.\\] Pick such an element $(gx,g,x)$. There exists $h\\in\\Gamma$ with $l(h)\\leq r_1$ and $hx\\in H^{(0)}$. Thus $(gx,gh^{-1},hx)$ and $(hx,h,x)$ are in $H^{+(r_1+r_2)}$ so $(gx,g,x)=(gx,gh^{-1},hx)(hx,h,x)\\in H^{+(r_1+r_2)}$.\n\\end{proof}\n\n\nNow we recall the definition of dynamical complexity for transformation groupoids. We refer the reader to \\cite[Appendix A]{GWY2} for a definition applicable to general \\'{e}tale groupoids.\n\n\\begin{defn}\\cite{GWY2} \\label{FDCbasic}\nLet $\\Gamma\\curvearrowright X$ be an action, let $G$ be an open subgroupoid of $\\Gamma\\ltimes X$, and let $\\mathcal{C}$ be a set of open subgroupoids of $\\Gamma\\ltimes X$. \nWe say that $G$ is decomposable over $\\mathcal{C}$ if for all $r\\geq 0$ there exists an open cover $G^{(0)}=U_0\\cup U_1$ of the unit space of $G$ such that for each $i\\in\\{0,1\\}$ the subgroupoid of $G$ generated by\n\\[ \\{ (gx,g,x)\\in G:x\\in U_i,l(g)\\leq r\\} \\]\nis in $\\mathcal{C}$.\n\nAn open subgroupoid of $\\Gamma\\ltimes X$ is said to have finite dynamical complexity if it is contained in the smallest class $\\mathcal{D}$ of open subgroupoids of $\\Gamma\\ltimes X$ that contains all relatively compact open subgroupoids and is closed under decomposability (i.e., if $G$ decomposes over $\\mathcal{D}$, then $G$ is in $\\mathcal{D}$).\n\nThe action $\\Gamma\\curvearrowright X$ is said to have finite dynamical complexity if $\\Gamma\\ltimes X$ has finite dynamical complexity.\n\\end{defn}\n\nThe following is a slight variation of the definition.\n\n\\begin{defn}\\cite{GWY2}\nLet $\\Gamma\\curvearrowright X$ be an action, let $G$ be an open subgroupoid of $\\Gamma\\ltimes X$, and let $\\mathcal{C}$ be a set of open subgroupoids of $\\Gamma\\ltimes X$. \nWe say that $G$ is strongly decomposable over $\\mathcal{C}$ if for all $r\\geq 0$ there exists an open cover $G^{(0)}=U_0\\cup U_1$ of the unit space of $G$ such that for each $i\\in\\{0,1\\}$, if $G_i$ is the subgroupoid of $G$ generated by\n\\[ \\{ (gx,g,x)\\in G:x\\in U_i,l(g)\\leq r\\}, \\]\nthen $G_i^{+r}$ (with expansion taken relative to $G$) is in $\\mathcal{C}$.\n\nDenote by $\\mathcal{D}_s$ the smallest class of open subgroupoids of $\\Gamma\\ltimes X$ that contains all relatively compact open subgroupoids and is closed under strong decomposability (i.e., if $G$ is strongly decomposable over $\\mathcal{D}_s$, then $G$ is in $\\mathcal{D}_s$).\n\\end{defn}\n\n\nThe following lemma describes properties of finite dynamical complexity that we will use later. We refer the reader to \\cite[Lemma 3.16]{GWY2} for the proof.\n\n\\begin{lem} \\label{FDCstrong} \nLet $\\Gamma\\curvearrowright X$ be an action.\n\\begin{enumerate}\n\\item If $G$ is an open subgroupoid of $\\Gamma\\ltimes X$ that is in $\\mathcal{D}$ (resp. $\\mathcal{D}_s$), then any open subgroupoid of $G$ is also in $\\mathcal{D}$ (resp. $\\mathcal{D}_s$).\n\\item $\\mathcal{D}=\\mathcal{D}_s$.\n\\end{enumerate}\n\\end{lem}\n\n\n\n\\section{Quantitative $K$-theory}\n\nIn this section, we recall some definitions and facts from our framework of quantitative (or controlled) $K$-theory in \\cite{Chung1}.\n\\begin{defn}\nA filtered Banach algebra is a Banach algebra $A$ with a family $(A_r)_{r\\geq 0}$ of linear subspaces such that\n\\begin{enumerate}\n\\item $A_{r_1}\\subset A_{r_2}$ if $r_1\\leq r_2$;\n\\item $A_{r_1}A_{r_2}\\subset A_{r_1+r_2}$ for all $r_1,r_2\\geq 0$;\n\\item $\\bigcup_{r\\geq 0}A_r$ is dense in $A$.\n\\end{enumerate}\nIf $A$ is unital with unit $1_A$, then we require $1_A\\in A_r$ for all $r\\geq 0$.\n\\end{defn}\n\nWe showed in Lemma \\ref{filtlem} that if $G$ is an open subgroupoid of $\\Gamma\\ltimes X$ and $s\\geq 0$, then $A^s(G)$ is a filtered Banach algebra with filtration \\[A^s(G)_r=\\{a\\in\\mathbb{C}_{L,0}[G;s]:\\mathrm{prop}_\\Gamma(a(t))\\leq r\\;\\text{for all}\\;t\\}.\\]\n\nWe will write $\\tilde{A}$ for the algebra obtained from $A$ by adjoining a unit if $A$ is non-unital, and $A$ itself if $A$ is already unital.\n\n\\begin{defn}\nLet $A$ be a filtered Banach algebra. For $0<\\varepsilon<\\frac{1}{20}$, $r\\geq 0$, and $N\\geq 1$,\n\\begin{enumerate}\n\\item an element $e\\in A$ is called an $(\\varepsilon,r,N)$-idempotent if $||e^2-e||<\\varepsilon$, $e\\in A_r$, and $\\max(||e||,||1_{\\tilde{A}}-e||)\\leq N$.\n\\item if $A$ is unital, an element $u\\in A$ is called an $(\\varepsilon,r,N)$-invertible if $u\\in A_r$, $||u||\\leq N$, and there exists $v\\in A_r$ with $||v||\\leq N$ such that $\\max(||uv-1||,||vu-1||)<\\varepsilon$.\n\\end{enumerate}\n\\end{defn}\n\n\\begin{defn}\nLet $A$ be a filtered Banach algebra.\n\\begin{enumerate}\n\\item Two $(\\varepsilon,r,N)$-idempotents $e_0$ and $e_1$ in $A$ are $(\\varepsilon',r',N')$-homotopic for some $\\varepsilon'\\geq\\varepsilon$, $r'\\geq r$, and $N'\\geq N$ if there exists a norm-continuous path $(e_t)_{t\\in[0,1]}$ of $(\\varepsilon',r',N')$-idempotents in $A$ from $e_0$ to $e_1$. \nEquivalently, there is an $(\\varepsilon',r',N')$-idempotent $e\\in C([0,1],A)$ such that $e(0)=e_0$ and $e(1)=e_1$.\n\\item If $A$ is unital, two $(\\varepsilon,r,N)$-invertibles $u_0$ and $u_1$ in $A$ are $(\\varepsilon',r',N')$-homotopic for some $\\varepsilon'\\geq\\varepsilon$, $r'\\geq r$, and $N'\\geq N$ if there exists an $(\\varepsilon',r',N')$-invertible $u$ in $C([0,1],A)$ with $u(0)=u_0$ and $u(1)=u_1$. In this case, we get a norm-continuous path $(u_t)_{t\\in[0,1]}$ of $(\\varepsilon',r',N')$-invertibles in $A$ from $u_0$ to $u_1$ by setting $u_t=u(t)$.\n\\end{enumerate}\n\\end{defn}\n\nGiven an $L^p$ operator algebra $A$ with a filtration $(A_r)_{r\\geq 0}$, one uses the above equivalence relations and goes through the standard procedure in the definition of $K$-theory of Banach algebras to arrive at the quantitative $K$-theory groups $K_*^{\\varepsilon,r,N}(A)$. We refer the reader to \\cite[Section 3]{Chung1} for details.\n\n\nThere are canonical homomorphisms \\[\\iota_*^{\\varepsilon,\\varepsilon',r,r',N,N'}:K_*^{\\varepsilon,r,N}(A)\\rightarrow K_*^{\\varepsilon',r',N'}(A)\\] for $0<\\varepsilon\\leq\\varepsilon'<\\frac{1}{20}$, $0\\leq r\\leq r'$, and $1\\leq N\\leq N'$, which we may think of as relaxation of control maps or inclusion maps.\n\nIf $e$ is an $(\\varepsilon,r,N)$-idempotent in a unital filtered Banach algebra $A$, then we may apply the holomorphic functional calculus to get an idempotent $c_0(e)\\in A$. This gives us a group homomorphism \\[c_0:K_0^{\\varepsilon,r,N}(A)\\rightarrow K_0(A).\\] Also, every $(\\varepsilon,r,N)$-invertible is actually invertible so we have a group homomorphism \\[c_1:K_1^{\\varepsilon,r,N}(A)\\rightarrow K_1(A)\\] given by $[u]_{\\varepsilon,r,N}\\mapsto[u]$. We sometimes refer to these homomorphisms as comparison maps.\n\n\n\\begin{prop}\\cite[Proposition 3.20]{Chung1} \\leavevmode \\label{qKtoKsurj}\n\\begin{enumerate}\n\\item Let $A$ be a filtered $L^p$ operator algebra. Let $f$ be an idempotent in $M_n(\\tilde{A})$, and let $0<\\varepsilon<\\frac{1}{20}$. Then there exist $r\\geq 0$ and $[e]\\in K_0^{\\varepsilon,r,||f||+1}(A)$ with $e$ an $(\\varepsilon,r,||f||+1)$-idempotent in $M_n(\\tilde{A})$ such that $c_0([e])=[f]$ in $K_0(A)$. \n\\item Let $A$ be a filtered $L^p$ operator algebra. Let $u$ be an invertible element in $M_n(\\tilde{A})$, and let $0<\\varepsilon<\\frac{1}{20}$. Then there exist $r\\geq 0$ and $[v]\\in K_1^{\\varepsilon,r,||u||+||u^{-1}||+1}(A)$ with $v$ an $(\\varepsilon,r,||u||+||u^{-1}||+1)$-invertible in $M_n(\\tilde{A})$ such that $c_1([v])=[u]$ in $K_1(A)$.\n\\end{enumerate}\n\\end{prop}\n\n\\begin{prop}\\cite[Proposition 3.21]{Chung1} \\leavevmode \\label{qKtoKinj}\n\\begin{enumerate}\n\\item There exists a non-decreasing function $P:[1,\\infty)\\rightarrow[1,\\infty)$ such that for any filtered $L^p$ operator algebra $A$, if $0<\\varepsilon<\\frac{1}{20P(N)}$, and $[e]_{\\varepsilon,r,N},[f]_{\\varepsilon,r,N}\\in K_0^{\\varepsilon,r,N}(A)$ are such that $c_0([e])=c_0([f])$ in $K_0(A)$, then there exist $r'\\geq r$ and $N'\\geq N$ such that in $K_0^{P(N)\\varepsilon,r',N'}(A)$ we have $[e]_{P(N)\\varepsilon,r',N'}=[f]_{P(N)\\varepsilon,r',N'}$.\n\\item Let $A$ be a filtered $L^p$ operator algebra. Suppose that $0<\\varepsilon<\\frac{1}{20}$, and $[u]_{\\varepsilon,r,N},[v]_{\\varepsilon,r,N}\\in K_1^{\\varepsilon,r,N}(A)$ are such that $c_1([u])=c_1([v])$ in $K_1(A)$. Then there exist $r'\\geq r$ and $N'\\geq N$ such that in $K_1^{\\varepsilon,r',N'}(A)$ we have $[u]_{\\varepsilon,r',N'}=[v]_{\\varepsilon,r',N'}$.\n\\end{enumerate}\n\\end{prop}\n\n\n\\begin{defn}\nA control pair is a pair $(\\lambda,h)$, where\n\\begin{enumerate}\n\\item $\\lambda:[1,\\infty)\\rightarrow[1,\\infty)$ is a non-decreasing function;\n\\item $h:(0,\\frac{1}{20})\\times[1,\\infty)\\rightarrow[1,\\infty)$ is a function such that $h(\\cdot,N)$ is non-increasing for fixed $N$. \\end{enumerate}\nWe will write $\\lambda_N$ for $\\lambda(N)$, and $h_{\\varepsilon,N}$ for $h(\\varepsilon,N)$.\n\nGiven two control pairs $(\\lambda,h)$ and $(\\lambda',h')$, we write $(\\lambda,h)\\leq(\\lambda',h')$ if $\\lambda_N\\leq \\lambda'_N$ and $h_{\\varepsilon,N}\\leq h'_{\\varepsilon,N}$ for all $\\varepsilon\\in(0,\\frac{1}{20})$ and $N\\geq 1$.\n\\end{defn}\n\n\nIn \\cite[Section 5]{Chung1}, we showed the existence of a controlled Mayer-Vietoris sequence for a pair of subalgebras satisfying certain conditions. Here, we state less general hypotheses (by considering pairs of ideals and omitting certain parameters) that suffice for our application. On the other hand, we also give ourselves a bit more flexibility in terms of propagation control. One can check that the proofs in \\cite{Chung1} carry over with appropriate adjustments of the propagation parameter $r$. Also note that our approach is slightly different from that in \\cite{GWY2}.\n\n\\begin{defn} \\label{MVpair}\nLet $A$ be a filtered $L^p$ operator algebra with filtration $(A_r)_{r\\geq 0}$. A pair $(I,J)$ of closed ideals of $A$ is a controlled Mayer-Vietoris pair for $A$ if it satisfies the following conditions:\n\\begin{enumerate}\n\\item There exists $\\rho:[0,\\infty)\\rightarrow[0,\\infty)$ with $\\rho(r)\\geq r$ such that for any $r\\geq 0$, any positive integer $n$, and any $x\\in M_n(A_r)$, there exist $x_1\\in M_n(I\\cap A_{\\rho(r)})$ and $x_2\\in M_n(J\\cap A_{\\rho(r)})$ such that $x=x_1+x_2$ and $\\max(||x_1||,||x_2||)\\leq ||x||$;\n\\item $I$ and $J$ have filtrations $(I\\cap A_r)_{r\\geq 0}$ and $(J\\cap A_r)_{r\\geq 0}$ respectively;\n\\item For any $r\\geq 0$, any $\\varepsilon>0$, any positive integer $n$, any $x\\in M_n(I\\cap A_r)$ and $y\\in M_n(J\\cap A_r)$ with $||x-y||<\\varepsilon$, there exists $z\\in M_n(I\\cap J\\cap A_{\\rho(r)})$ such that $\\max(||z-x||,||z-y||)<\\varepsilon$, where $\\rho$ is as above.\n\\end{enumerate}\n\\end{defn}\n\n\n\n\\begin{thm}\\cite[Definition 5.12 and Theorem 5.14]{Chung1} \\label{MVthm}\nThere exist control pairs $(\\lambda,h)\\leq(\\lambda',h')$ such that for any filtered $L^p$ operator algebra $A$ and any controlled Mayer-Vietoris pair $(I,J)$ for $A$, if $x\\in K_1^{\\varepsilon,r,N}(A)$, then there exists \\[\\partial_c(x)\\in K_0^{\\lambda_N\\varepsilon,h_{\\varepsilon,N}r,\\lambda_N}(I\\cap J)\\] such that if $\\partial_c(x)=0$, then there exist $a\\in K_0^{\\lambda'_N\\varepsilon,h'_{\\varepsilon,N}r,\\lambda'_N}(I)$ and $b\\in K_0^{\\lambda'_N\\varepsilon,h'_{\\varepsilon,N}r,\\lambda'_N}(J)$ such that $x=a+b$ in $K_0^{\\lambda'_N\\varepsilon,h'_{\\varepsilon,N}r,\\lambda'_N}(A)$.\n\\end{thm}\n\n\n\n\n\n\n\n\n\\section{Main theorem}\n\nIn this section, we shall prove the main theorem that the $L^p$ assembly map for $\\Gamma\\curvearrowright X$ in Definition \\ref{assembly} is an isomorphism if $\\Gamma\\curvearrowright X$ has finite dynamical complexity.\nThe proof that we present is modeled after the proof in the $C^*$-algebraic setting in \\cite{GWY2}, consisting of a homotopy invariance argument and a Mayer-Vietoris argument.\n\n\n\n\n\n\n\n\\subsection{Homotopy invariance} \\label{sec:Base}\n\nRecall that we use the shorthand $A^s(G)$ for $B^p_{L,0}(G;s)$. The homotopy invariance argument involved in the proof of the main theorem is given by the following statement. \n\n\\begin{prop} \\label{baseprop}\nLet $G$ be an open subgroupoid of $\\Gamma\\ltimes X$ such that \\[G\\subseteq\\{(gx,g,x)\\in\\Gamma\\ltimes X:l(g)\\leq s\\}\\] for some $s\\geq 0$. Then $K_*(A^s(G))=0$.\n\\end{prop}\n\n\n\nBefore getting into the proof of the proposition, we need to fix some terminology that is standard in the $C^*$-algebraic setting but perhaps less so when Hilbert spaces are replaced by other Banach spaces. Having done that, the series of results in the rest of this section then yields the proposition.\n\n\n\\begin{defn}\nLet $E$ be a complex Banach space. We say that $T\\in B(E)$ is a partial isometry if $||T||\\leq 1$ and there exists $S\\in B(E)$ such that $||S||\\leq 1$, $TST=T$, and $STS=S$. We call such an $S$ a generalized inverse of $T$.\n\\end{defn}\n\n\\begin{rem}\t\\leavevmode\n\\begin{enumerate}\n\\item In \\cite[Section 6]{Phil12}, Phillips considers spatial partial isometries on $L^p$ spaces. Such spatial partial isometries are partial isometries in the sense of the preceding definition but the converse is not true.\n\\item If $(Z,\\mu)$ is a $\\sigma$-finite measure space, $p\\in[1,\\infty)\\setminus\\{2\\}$, and $T\\in B(L^p(Z,\\mu))$ is an isometric (but not necessarily surjective) linear map, then it follows from Lamperti's theorem \\cite{Lamp} (also see \\cite[Theorem 6.9]{Phil12}) that $T$ is a partial isometry in the sense above, and one can find a generalized inverse $S$ such that $ST=I$ (cf. \\cite[Section 6]{Phil12}). Hereafter, we will denote such an $S$ by $T^\\dagger$.\n\\end{enumerate}\n\\end{rem}\n\n\n\n\\begin{defn}\nIf $A\\subset B(L^p(\\mu))$ is an $L^p$ operator algebra, then we say that $b\\in B(L^p(\\mu))$ is a multiplier of $A$ if $bA\\subset A$ and $Ab\\subset A$. We say that $b$ is an isometric multiplier of $A$ if $b$ is an isometry and both $b$ and $b^\\dagger $ are multipliers of $A$. Denote by $M(A)$ the set of all multipliers of $A$.\n\\end{defn}\n\nNote that $M(A)$ is also an $L^p$ operator algebra.\n\n\n\n\nFor $G$ an open subgroupoid of $\\Gamma\\ltimes X$ and $s\\geq 0$, we introduce the following notation: \n\\begin{itemize}\n\\item $P_s(G)=\\{(z,x)\\in P_s(\\Gamma)\\times X:(gx,g,x)\\in G\\;\\text{for all}\\;g\\in \\mathrm{supp}(z)\\}$.\n\\item $Z_G=(Z_s\\times X)\\cap P_s(G)$.\n\\item $E_G=\\ell^p(Z_G,\\ell^p\\otimes\\ell^p(\\Gamma))\\cong\\ell^p(Z_G)\\otimes\\ell^p\\otimes\\ell^p(\\Gamma)$.\n\\end{itemize}\n\nNote that $E_G$ is a subspace of $E_s$. \nMoreover, the faithful representation of $B^p(G;s)$ on $E_s$ restricts to a faithful representation on $E_G$. \nThus we will regard $B^p(G;s)$ as faithfully represented on $E_G$, and $A^s(G):=B^p_{L,0}(G;s)$ as faithfully represented on $L^p([0,\\infty),E_G)$.\n\nIf $(z,x)\\in P_s(G)$ and $\\mathrm{supp}(z)=\\{g_1,\\ldots,g_n\\}$, then $\\{e,g_1,\\ldots,g_n\\}$ also spans a simplex $\\Delta$ in $P_s(\\Gamma)$ such that $\\Delta\\times\\{x\\}$ is contained in $P_s(G)$. Hence the family of functions \\begin{equation}\\label{homodef} F_r:P_s(G)\\rightarrow P_s(G), (z,x)\\mapsto ((1-r)z+re,x) \\quad\\quad (0\\leq r\\leq 1) \\end{equation} defines a homotopy between the identity map on $P_s(G)$ and the projection onto the subset $\\{(z,x)\\in P_s(G):z=e\\}$, which we may identify with the unit space $G^{(0)}$.\n\nIn the definition of $\\mathbb{C}[\\Gamma\\curvearrowright X;s]$, we may use $\\mathcal{K}_\\Gamma^\\infty$, the algebra of compact operators on $(\\bigoplus_{n=0}^\\infty \\ell^p\\otimes\\ell^p(\\Gamma))_{p}\\cong (\\bigoplus_{n=0}^\\infty\\ell^p)_{p}\\otimes\\ell^p(\\Gamma)$, thereby obtaining another $L^p$ Roe algebra $B^p(\\Gamma\\curvearrowright X;\\mathcal{K}_\\Gamma^\\infty;s)$. Moreover, fixing an isometric isomorphism $\\phi:\\ell^p\\stackrel{\\cong}{\\rightarrow}(\\bigoplus_{n=0}^\\infty\\ell^p)_{p}$ gives an isomorphism \\[B^p(\\Gamma\\curvearrowright X;s)\\cong B^p(\\Gamma\\curvearrowright X;\\mathcal{K}_\\Gamma^\\infty;s).\\]\nWe also have the corresponding statements for the $L^p$ localization algebras and obstruction algebras defined earlier. \n\nFor each $n$, define an isometry $u_{n,0}:\\ell^p\\rightarrow(\\bigoplus_{n=0}^\\infty\\ell^p)_{p}$ by inclusion as the $n$th summand, and define $u_{n,0}^\\dagger :(\\bigoplus_{n=0}^\\infty\\ell^p)_{p}\\rightarrow\\ell^p$ by projection onto the $n$th summand. Then $u_{n,0}^\\dagger u_{n,0}=I$ for all $n$, and $u_{n,0}^\\dagger u_{m,0}=0$ when $n\\neq m$. Define $u_n:L^p([0,\\infty),E_G)\\rightarrow L^p([0,\\infty),E_G^\\infty)$ to be the operator induced by tensoring $u_{n,0}$ with the identity on the other factors, where $E_G^\\infty=(\\bigoplus_{n=0}^\\infty E_G)_{p}$, and define $u_n^\\dagger $ similarly using $u_{n,0}^\\dagger $. Then $u_n^\\dagger u_n=I$ for all $n$, and $u_n^\\dagger u_m=0$ when $n\\neq m$. \n\nGiven $a\\in B(\\ell^p)$, consider $a^\\infty=a\\oplus a\\oplus\\cdots\\in B((\\bigoplus_{n=0}^\\infty\\ell^p)_{p})$. Then $\\mu(a)=a^\\infty$ is an isometric homomorphism $B(\\ell^p)\\rightarrow B((\\bigoplus_{n=0}^\\infty\\ell^p)_{p})$. We may also consider the isometry $v\\in B((\\bigoplus_{n=0}^\\infty\\ell^p)_{p})$ given by the right shift taking the $n$th summand onto the $(n+1)$st summand. Denote by $v^\\dagger \\in B((\\bigoplus_{n=0}^\\infty\\ell^p)_{p})$ the left shift. Then $v\\mu(a)v^\\dagger =0\\oplus a\\oplus a\\oplus\\cdots$ for all $a\\in B(\\ell^p)$. With $u_{0,0}$ as above, $u_{0,0}au_{0,0}^\\dagger $ is given by $a\\oplus 0\\oplus 0\\oplus\\cdots$ for all $a\\in B(\\ell^p)$. Now $a\\mapsto\\mu^{+1}(a):=v\\mu(a)v^\\dagger $ and $a\\mapsto\\mu^0(a):=u_{0,0}au_{0,0}^\\dagger $ are bounded homomorphisms $B(\\ell^p)\\rightarrow B((\\bigoplus_{n=0}^\\infty\\ell^p)_{p})$.\n\nNow given $a\\in B(L^p([0,\\infty),E_G))$, consider the bounded linear operator $\\mu(a)=a^\\infty$ on $L^p([0,\\infty),E_G^\\infty)$. Proceeding similarly as above, we get bounded homomorphisms \\[\\mu,\\mu^{+1},\\mu^0:B(L^p([0,\\infty),E_G))\\rightarrow B(L^p([0,\\infty),E_G^\\infty)).\\] Moreover, each of them maps $A^s(G)$ into $A^s(G;\\mathcal{K}_\\Gamma^\\infty)$, and $M(A^s(G))$ into $M(A^s(G;\\mathcal{K}_\\Gamma^\\infty))$.\n\n\nThe following lemma is an $L^p$ version of a fairly standard result in the $K$-theory of $C^*$-algebras and can be proved in the same way as it is done for $C^*$-algebras (cf. \\cite[Lemma 4.6.2]{HR} or \\cite[Section 3, Lemma 2]{HRY}).\n\\begin{lem} \\label{lemMult}\nLet $\\alpha:A\\rightarrow C$ be a bounded homomorphism of $L^p$ operator algebras with $C\\subset B(L^p(\\mu))$, and let $v\\in B(L^p(\\mu))$ be an isometric multiplier of $C$. Then the map $a\\mapsto v\\alpha(a)v^\\dagger $ is a bounded homomorphism from $A$ to $C$, and induces the same map as $\\alpha$ on $K$-theory.\n\nMore generally, if $v\\in B(L^p(\\mu))$ is a partial isometry and a multiplier of $C$, $w$ is a generalized inverse of $v$ that is also a multiplier of $C$, and $\\alpha(a)wv=\\alpha(a)=wv\\alpha(a)$ for all $a\\in A$, then the map $a\\mapsto v\\alpha(a)w$ is a bounded homomorphism from $A$ to $C$, and induces the same map as $\\alpha$ on $K$-theory.\n\\end{lem}\n\n\\begin{lem}\n$K_*(M(A^s(G)))=0$.\n\\end{lem}\n\n\\begin{proof}\nNote that $\\mu,\\mu^{+1}:M(A^s(G))\\rightarrow M(A^s(G;\\mathcal{K}_\\Gamma^\\infty))$ induce the same map on $K$-theory by the previous lemma. Moreover, since $\\mu^0(a)\\mu^{+1}(a)=\\mu^{+1}(a)\\mu^0(a)=0$ for all $a\\in M(A^s(G))$ and $\\mu=\\mu^0+\\mu^{+1}$, the induced maps on $K$-theory satisfy $\\mu_*=\\mu^0_*+\\mu^{+1}_*=\\mu^0_*+\\mu_*$. Hence $\\mu^0_*=0$. But $\\mu^0$ induces an isomorphism on $K$-theory so $K_*(M(A^s(G)))=0$.\n\\end{proof}\n\n\nFor each $z\\in Z_s$ such that $(z,x)\\in P_s(G)$ for some $x\\in X$, let $E_z$ be a copy of $\\ell^p$ so that we have an isometric isomorphism $\\ell^p\\cong(\\bigoplus_{z\\in Z}E_z)_{p}$, and let $w_z:\\ell^p\\otimes\\ell^p(\\Gamma)\\rightarrow \\ell^p\\otimes\\ell^p(\\Gamma)$ be an isometry with range $E_z\\otimes\\ell^p(\\Gamma)$. For each $r\\in\\mathbb{Q}\\cap[0,1]$, define an isometry \\[w(r):\\ell^p(Z_G)\\otimes\\ell^p\\otimes\\ell^p(\\Gamma)\\rightarrow\\ell^p(Z_G)\\otimes\\ell^p\\otimes\\ell^p(\\Gamma)\\] by $\\delta_{z,x}\\otimes\\eta\\mapsto\\delta_{(1-r)z+re,x}\\otimes w_z\\eta$. Note that if $r\\neq s$, then $w(r)^\\dagger w(s)=0$.\n\nFor $t\\geq 0$ and $n\\in\\mathbb{N}\\cup\\{\\infty\\}$, define an isometry \\[v_n(t):\\ell^p(Z_G)\\otimes\\ell^p\\otimes\\ell^p(\\Gamma)\\rightarrow\\ell^p(Z_G)\\otimes\\ell^p\\otimes\\ell^p(\\Gamma)\\] in the following way.\nFor $m\\in\\mathbb{N}$, define \\[ v_n(m)=\\begin{cases} w(0) & \\text{if}\\;m\\leq n \\\\ w(\\frac{m-n}{n}) & \\text{if}\\;n2m$. It follows that the Rips-propagation of $v_n(t)a(t)v_n(t)^\\dagger $ is at most $\\mathrm{prop}_{Rips}a(t)+\\min(1,\\frac{2}{|t-1|})$, so $v_nav_n^\\dagger \\in\\mathbb{C}_{L,0}[G;s]$.\n\nAlso, the operators $S_t:=v_{n+1}(t)v_n(t)^\\dagger $ on $\\ell^p(Z_G)\\otimes\\ell^p\\otimes\\ell^p(\\Gamma)$ have matrix entries $(S_t)_{y,z}$ that act as constant functions $X\\rightarrow B(\\ell^p\\otimes\\ell^p(\\Gamma))$, their Rips-propagation tends to zero as $t\\rightarrow\\infty$, and they have $\\Gamma$-propagation at most $s$ for all $t$. Hence $v_{n+1}v_n^\\dagger $ is a multiplier of $A^s(G)$ for all $n$.\n\n\\begin{lem}\nLet $A$ be a unital Banach algebra, and let $I$ be an ideal in $A$. Define the double of $A$ along $I$ to be $D=\\{(a,b)\\in A\\oplus A:a-b\\in I\\}$. Assume that $A$ has trivial $K$-theory. Then the inclusion $\\iota:I\\rightarrow D$ given by $a\\mapsto(a,0)$ induces an isomorphism in $K$-theory, and the diagonal inclusion $\\delta:I\\rightarrow D$ given by $a\\mapsto(a,a)$ induces the zero map on $K$-theory.\n\\end{lem}\n\n\\begin{proof}\nNote that $\\iota(I)$ is an ideal in $D$, and $D\/\\iota(I)$ is isomorphic to $A$ via the second coordinate projection. Since $K_*(A)=0$, it follows from the six-term exact sequence that $\\iota$ induces an isomorphism in $K$-theory. \n\nOn the other hand, $\\delta$ factors through the diagonal inclusion $A\\rightarrow D,a\\mapsto(a,a)$, so $\\delta$ induces the zero map on $K$-theory since $K_*(A)=0$.\n\\end{proof}\n\nWe shall apply the lemma in the case where $A=M(A^s(G))$ and $I=A^s(G)$ to prove the next proposition.\n\n\\begin{prop}\nLet $v_n:L^p([0,\\infty),E_G)\\rightarrow L^p([0,\\infty),E_G)$ be as defined above. Then the maps $a\\mapsto v_0 av_0^\\dagger $ and $a\\mapsto v_\\infty av_\\infty^\\dagger $ induce the same map $K_*(A^s(G))\\rightarrow K_*(A^s(G))$.\n\\end{prop}\n\n\\begin{proof}\nGiven $a\\in A^s(G)$, define \\[\\alpha(a)=(\\bigoplus_{n=0}^\\infty v_nav_n^\\dagger ,\\bigoplus_{n=0}^\\infty v_\\infty av_\\infty^\\dagger )\\in A^s(G;\\mathcal{K}_\\Gamma^\\infty)\\oplus A^s(G;\\mathcal{K}_\\Gamma^\\infty).\\] Also define \\[\\beta(a)=(\\bigoplus_{n=0}^\\infty v_{n+1}av_{n+1}^\\dagger ,\\bigoplus_{n=0}^\\infty v_\\infty av_\\infty^\\dagger )\\in A^s(G;\\mathcal{K}_\\Gamma^\\infty)\\oplus A^s(G;\\mathcal{K}_\\Gamma^\\infty).\\] Let $D$ be the double of $M(A^s(G;\\mathcal{K}_\\Gamma^\\infty))$ along $A^s(G;\\mathcal{K}_\\Gamma^\\infty)$, and let \\[C=\\{(c,d)\\in D:d=\\bigoplus_{n=0}^\\infty v_\\infty av_\\infty^\\dagger \\;\\text{for some}\\;a\\in A^s(G)\\},\\] which is a closed subalgebra of $D$. Moreover, $\\alpha$ and $\\beta$ are bounded homomorphisms with image in $C$. \nConsider $w=(w_1,w_2)$, where $w_1=\\bigoplus_{n=0}^\\infty v_{n+1}v_n^\\dagger $ and $w_2=\\bigoplus_{n=0}^\\infty v_\\infty v_\\infty^\\dagger $. Note that $w_1,w_2\\in M(A^s(G;\\mathcal{K}_\\Gamma^\\infty))$. We claim that $w$ is a multiplier of $C$. Indeed, if $(c,d)\\in C$, then $w_2d=dw_2=d$ so it suffices to show that $cw_1-d$ and $w_1c-d$ are in $A^s(G;\\mathcal{K}_\\Gamma^\\infty)$. We will only consider $w_1c-d$ since the other case is similar. Now $w_1c-d=w_1(c-d)+(w_1d-d)$ so it suffices to show that $w_1d-d\\in A^s(G;\\mathcal{K}_\\Gamma^\\infty)$. But $w_1d-d=(w_1-w_2)d=\\bigoplus_{n=0}^\\infty(v_{n+1}v_n^\\dagger -v_\\infty v_\\infty^\\dagger )v_\\infty a v_\\infty^\\dagger \\in A^s(G;\\mathcal{K}_\\Gamma^\\infty)$ since $v_n(t)=v_\\infty(t)$ for each fixed $t$ and all $n>t$. Similarly, $w^\\dagger =(w_1^\\dagger ,w_2^\\dagger )$ is a multiplier of $C$. \n\nNow $\\beta(a)=w\\alpha(a)w^\\dagger $ for all $a\\in A$. Moreover, $\\alpha(a)w^\\dagger w=\\alpha(a)=w^\\dagger w\\alpha(a)$ for all $a\\in A$ so $\\alpha$ and $\\beta$ induce the same map $K_*(A^s(G))\\rightarrow K_*(C)$ by Lemma \\ref{lemMult}, and thus the same map $K_*(A^s(G))\\rightarrow K_*(D)$ upon composing with the map induced by the inclusion of $C$ into $D$.\n\nLet $u$ be the isometric multiplier of $A^s(G;\\mathcal{K}_\\Gamma^\\infty)$ induced by the right shift. Then $(u,u)$ is a multiplier of $D$, and conjugating $\\beta(a)$ by $(u,u)$ gives \\[\\gamma(a)=\\biggl(0\\oplus\\bigoplus_{n=1}^\\infty v_nav_n^\\dagger ,0\\oplus\\bigoplus_{n=1}^\\infty v_\\infty av_\\infty^\\dagger \\biggr).\\] Thus $\\beta$ and $\\gamma$ induce the same map $K_*(A)\\rightarrow K_*(D)$. On the other hand, the homomorphism $\\delta:A^s(G)\\rightarrow D$ given by \\[a\\mapsto(v_\\infty av_\\infty^\\dagger \\oplus 0\\oplus 0\\oplus\\cdots,v_\\infty av_\\infty^\\dagger \\oplus 0\\oplus 0\\oplus\\cdots)\\] induces the zero map on $K$-theory by the previous lemma. Also, $\\gamma(a)\\delta(a)=\\delta(a)\\gamma(a)=0$. Hence \\[\\alpha_*=\\beta_*=\\gamma_*=\\gamma_*+\\delta_*=(\\gamma+\\delta)_*:K_*(A^s(G))\\rightarrow K_*(D).\\]\n\nLet $\\psi_0,\\psi_\\infty:A^s(G)\\rightarrow D$ be the homomorphisms defined by \\begin{align*} \\psi_0(a) &=(v_0 av_0^\\dagger \\oplus 0\\oplus 0\\oplus\\cdots,0), \\\\ \\psi_\\infty(a) &=(v_\\infty av_\\infty^\\dagger \\oplus 0\\oplus 0\\oplus\\cdots,0).\\end{align*} Also define $\\zeta:A^s(G)\\rightarrow D$ by \\[\\zeta(a)=\\biggl(0\\oplus\\bigoplus_{n=1}^\\infty v_nav_n^\\dagger ,\\bigoplus_{n=0}^\\infty v_\\infty av_\\infty^\\dagger \\biggr).\\] Note that $\\zeta(a)\\psi_0(a)=\\psi_0(a)\\zeta(a)=\\zeta(a)\\psi_\\infty(a)=\\psi_\\infty(a)\\zeta(a)=0$ for all $a\\in A^s(G)$. Also, $\\psi_0+\\zeta=\\alpha$ and $\\psi_\\infty+\\zeta=\\gamma+\\delta$. Hence \\[(\\psi_0)_*+\\zeta_*=\\alpha_*=(\\gamma+\\delta)_*=(\\psi_\\infty)_*+\\zeta_*,\\] so $\\psi_0$ and $\\psi_\\infty$ induce the same maps on $K$-theory.\n\nFinally, if $\\iota:A^s(G)\\rightarrow D$ is the inclusion into the first factor (where $D$ is now regarded as the double of $M(A^s(G))$ along $A^s(G)$), then $\\psi_i(a)$ is given by the composition\n\\[ a\\mapsto v_iav_i^\\dagger \\stackrel{\\iota}{\\mapsto}(v_iav_i^\\dagger ,0)\\mapsto (v_iav_i^\\dagger \\oplus 0\\oplus 0\\oplus\\cdots,0), \\] and the last two maps induce isomorphisms on $K$-theory, so $a\\mapsto v_0av_0^\\dagger $ and $a\\mapsto v_\\infty av_\\infty^\\dagger$ induce the same map on $K$-theory.\n\\end{proof}\n\n\nNow it remains to be shown that $a\\mapsto v_\\infty av_\\infty^\\dagger $ induces the identity map on $K_*(A^s(G))$ while $a\\mapsto v_0 av_0^\\dagger $ induces the zero map on $K_*(A^s(G))$. This will complete the proof of Proposition \\ref{baseprop}.\n\n\\begin{lem}\nThe map $\\phi_\\infty:K_*(A^s(G))\\rightarrow K_*(A^s(G))$ induced by conjugation by $v_\\infty$ is the identity map.\n\\end{lem}\n\n\\begin{proof}\nSince $v_\\infty(t)=w(0)$ for all $t$, and $w(0)$ is an isometric multiplier of $A^s(G)$, $\\phi_\\infty$ induces the identity map on $K$-theory by Lemma \\ref{lemMult}.\n\\end{proof}\n\n\\begin{lem}\nThe map $\\phi_0:K_*(A^s(G))\\rightarrow K_*(A^s(G))$ induced by conjugation by $v_0$ is the zero map.\n\\end{lem}\n\n\\begin{proof}\nLet $G^{(0)}$ be the unit space of $G$, which is an open subgroupoid of $\\Gamma\\ltimes X$. We may then consider $A^s(G^{(0)})$. In fact, $\\phi_0$ factors through $K_*(A^s(G^{(0)}))$, i.e., we have a commutative diagram\n\\[\n\\begindc{\\commdiag}[100]\t\n\\obj(10,5)[2a]{$K_*(A^s(G))$}\t\\obj(20,5)[2b]{$K_*(A^s(G))$}\t\n\n\\mor{2a}{2b}{$\\phi_0$}\t \n\n\\obj(15,0)[3a]{$K_*(A^s(G^{(0)}))$}\t\n\n\\mor{2a}{3a}{}\t\\mor{3a}{2b}{} \t\n\\enddc\n\\]\nso it suffices to show that $K_*(A^s(G^{(0)}))=0$.\n\nNote that $a(t)$ has zero Rips-propagation for all $a\\in A^s(G^{(0)})$ and $t\\in[0,\\infty)$. For each $n\\in\\mathbb{N}$ and $a\\in A^s(G^{(0)})$, define \\[a^{(n)}(t)=\\begin{cases} a(t-n) & \\text{for}\\; t\\geq n \\\\ 0 & \\text{for}\\; t