diff --git a/data_all_eng_slimpj/shuffled/split2/finalzfqb b/data_all_eng_slimpj/shuffled/split2/finalzfqb new file mode 100644 index 0000000000000000000000000000000000000000..0e6a51813dbb4a4bb061dac8123fb1c7eddf7691 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzfqb @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\nIn sexual reproduction, as opposed to asexual reproduction, the genomes of the \ntwo parents are mixed, and within the diploid genome of each parent happens\ncrossover. This way of reproduction has advantages as well as disadvantages \ncompared with asexual cloning of haploid genomes. An advantage is that bad \nrecessive mutations do not affect the health if they are present in only one of \nthe two haplotypes (= sets of genetic information). A disadvantage\nis the reduced number of births if only the females produce offspring while the \nmales consume as much food and space as the females. Moreover, crossover of \ntwo different genomes may produce a mixture which is fitter than each the two \nparents but also one which is less fit, as seen these \ndays in the DaimlerChrysler car company (outbreeding depression). For small \npopulations, the probability is higher that the two parents have the same bad \nrecessive mutation which therefore diminishes the health of the individual \n(inbreeding depression). \n\n\\begin{figure}[hbt]\n\\begin{center}\n\\includegraphics[angle=-90,scale=0.5]{inbreeding1.ps}\n\\end{center}\n\\caption{Average over 100 simulations with $G=10$ groups each, $\\Delta = 100$.\n}\n\\end{figure}\n\n\\begin{figure}[hbt]\n\\begin{center}\n\\includegraphics[angle=-90,scale=0.5]{inbreeding2.ps}\n\\end{center}\n\\caption{Average over 10 simulations with a large (top) and 100 with a small\n(bottom) population, versus number $G$ of groups; $r=1, \\; \\Delta = 1000$ and \n100, respectively. For the larger population $G$ is divided by 10 so that \ndata for the same number of individuals per group have the same horizontal\ncoordinate. We see nice scaling.\n}\n\\end{figure}\n\n\\begin{figure}[hbt]\n\\begin{center}\n\\includegraphics[angle=-90,scale=0.30]{inbreeding3a.ps}\n\\includegraphics[angle=-90,scale=0.30]{inbreeding3b.ps}\n\\end{center}\n\\caption{Time dependence of outbreeding advantage (part a) and outbreeding\ndepression (part b).\n}\n\\end{figure}\n\n\\begin{figure}[hbt]\n\\begin{center}\n\\includegraphics[angle=-90,scale=0.5]{inbreeding4.ps}\n\\end{center}\n\\caption{Average over 100 simulations with a small population, high minimum\nage of reproduction, and $G=1$ and 50. For $G=1$ there is always complete \nmixing. Note the double-logarithmic scales, also in Fig.5 and 6.\n}\n\\end{figure}\n\n\\begin{figure}[hbt]\n\\begin{center}\n\\includegraphics[angle=-90,scale=0.5]{inbreeding5.ps}\n\\end{center}\n\\caption{Average over 100 simulations with a small population, high minimum\nage of reproduction and various\nlengths $L$ of the bit-strings, using a birth rate $B = 128\/L$. \n}\n\\end{figure}\n\n\\begin{figure}[hbt]\n\\begin{center}\n\\includegraphics[angle=-90,scale=0.5]{inbreeding6.ps}\n\\end{center}\n\\caption{Dependence on population size for $N_{\\max} = 10^3 \\dots 10^7$, \naveraged over 1000 to one sample. $L=64, \\; B = 2$.\n}\n\\end{figure}\n\n\\begin{figure}[hbt]\n\\begin{center}\n\\includegraphics[angle=-90,scale=0.5]{inbreeding11.ps}\n\\end{center}\n\\caption{\nRelation between normalized population size and the crossover\nrate. The population size was divided by the population evolving with crossover\nrate $r = 1$.\n}\n\\end{figure}\n\n\\begin{figure}[hbt]\n\\begin{center}\n\\includegraphics[angle=-90,scale=0.5]{inbreeding12.ps}\n\\end{center}\n\\caption{\nDistribution of defective genes in the genomes of populations\nevolving under different crossover frequencies.\n}\n\\end{figure}\n\n\\begin{figure}[hbt]\n\\begin{center}\n\\includegraphics[angle=-90,scale=0.5]{inbreeding13.ps}\n\\end{center}\n\\caption{\nRelation between crossover frequency and average frequency of\ndefective genes in the sections of genomes expressed before the\nreproduction age $R$. Note, fractions equal 0 mean that the populations with \nthe given crossover frequency died out.\n}\n\\end{figure}\n\n\\section{Standard Model}\n\nWe try to simulate these effects in the standard sexual Penna ageing model,\ndeviating from published programs \\cite{books} as follows: The Verhulst factor, \na death probability $N\/N_{\\max}$ at population $N$ with carrying capacity\n$N_{\\max} = 10^3, \\ 10^4, \\,10^5$,\ndue to limited food and space, was applied to the births only and not to adults;\nthe initial age distribution was taken as random when needed;\nthe birth rate was reduced from 4 to 1, the lethal threshold of active \nmutations from 3 to 1 (that means a single active mutation kills), and mostly \nonly $10^4$ instead of $2 \\times 10^4$ time steps were made. (One time step or\niteration is one Monte Carlo step for each individual). Furthermore\nthe whole population was for most of the simulated time separated into $G$ \ndifferent groups such that females look for male partners only within their\nown group, with a separate Verhulst factor applying to each group. For the\nlast $\\Delta \\ll 10^4$ time steps this separation into groups was dissolved: \nThen females could select any male, and only one overall Verhulst factor \napplied to the whole population. Finally, the crossover\nprocess within each parent before each birth was not made always but only\nwith a crossover probability $r$. \n\nIf there would be no inbreeding depression then during the first longer part\nof the simulation the total number $N_1$ of individuals would be independent of \nthe number $G$ of groups into which it is divided. And if then there are no \nadvantages or disadvantages of outbreeding, the population $N_2$ during the \nsecond shorter part, $10^4 - \\Delta < t < 10^4$, would be the same as the\npreceding population $N_1$ during the last section, $10^4 - 2 \\Delta < t < \n10^4 - \\Delta$, of the longer first part. We will present data showing that \nthis is not the case. Similar simulations for $G=2$ groups was published\nlong ago \\cite{LSC}. \n \nA difficulty in such simulations is the Eve effect: After a time \nproportional to the population size, everybody has the same female (Eve) and \nthe same male (Adam) as ancestor, with all other offspring having gotten less \nfit genomes due to random mutations and thus having died out. If we would\ndivide the whole population into many groups without further changes, the Eve \neffect would let all groups but one die out and thus destroy the separation.\nTherefore for the first long period of separation we used separate Verhulst \nfactors for each group, stabilizing its population, while for the second shorter\npart of mixing we used mostly $\\Delta = N_{\\max}\/100$. \n\nFigure 1 shows the dependence on the crossover probability for the \npopulations $N_1$ before and $N_2$ after mixing. We see that the mixing always\nincreases the population, that means one has no outbreeding depression but\nan outbreeding advantage. Figure 2 confirms this advantage but also \nshows the inbreeding depression: The larger the number $G$ of groups (and \nthus the smaller the group size) is, the smaller are the two populations\n$N_1$ and $N_2$. (The difference between $N_1$ and $N_2$ fluctuates less than\nthese numbers themselves since $N_2$ is strongly correlated with $N_1$.)\nAlso, for the larger population in Fig.2, the number of groups can be larger\nbefore the population becomes extinct. Figure 3 shows the time dependence \nof the outbreeding effect with mixing between groups allowed after 9900 (part a)\nand 9000 (part b) time steps. Figure 3a shows summed populations from \n100 simulations with a small population ($G = 10$) and 10 simulations of a large\npopulation ($G = 100$), versus time after mixing started; $r = 1$ in both cases.\nFor much larger populations of 5 million and still $G = 10$, no such effect of \nmixing is seen. Part b shows for the high\nreproduction age $R$ of the following figures one example of the outbreeding\ndepression (bottleneck \\cite{malarz}) followed by a recovery with oscillations \nof period $R$ after mixing was allowed from time 9001 on; $N_{\\max}=10^5$..\n\n \nWe also checked for the influence of $r$ in the case when the minimum age of \nreproduction $R$ is 5\/8 of the length $L$ of the bit-strings, i.e. larger than \nthe value of 8 used before, and when $L$ is different from the 32 used in Figs.\n1 to 3. In these simulations we also assumed all mutations to be\nrecessive, in contrast to the 6 out of 32 dominant bit positions for Figs.\n1 to 3. Figure 4 shows for $L = 32$ and a birth rate $B = 4$ \na minimum of the population at intermediate $r$ for one group, and for 50 \ngroups a monotonic behaviour but with outbreeding depression at small $r$ and\noutbreeding advantage at big $r$. This population minimum is seen for\n$L = 64$ and 32 but not for 16, Fig.5. Figure 6 shows the dependence on \npopulation size. (Our data before and after mixing are average over $\\Delta = \n100$ or 1000 iterations. When outbreeding depression occurs it may happen that \nlater the population recovers: Fig.3b.)\n\n\\section{Interpretation}\n\nTo study the inbreeding and outbreeding depressions in detail we have\nanalyzed the results of simulations of single populations of different size\nunder different regime of intragenomic recombinations (crossover rate $r$).\nParameters for these simulations have been slightly changed to get clearer \nresults: $L= 128,\\; R = 80,\\; N_{\\max} = 1000$ to 20 000, crossover probability \n$r = 0$ to 1, $B = 1$, time of simulations = $5 \\times 10^5$ iterations.\nIn Fig. 7 the relation between the size of population and the crossover\nprobability for three different environment capacities are shown.\n\nPopulations in the smallest environment ($N_{\\max}=1000$) survive with $r = 0$ \nbut their sizes decrease with increasing $r$ and are extinct for $r$ set between\n0.12 and 0.4. Under larger crossover rates populations survive and their\nsizes are larger than those obtained for $r = 0$ (see plots in Fig. 7 where\nsizes of populations were normalized by the size of population under \n$r=1$). Larger populations ($N_{\\max}= 10 000$) are extinct in a very\nnarrow range of crossover rates close to 0.12, and populations with $N_{\\max}= \n20000$ become extinct at slightly lower crossover rates. Nevertheless, all\npopulations have larger sizes when the crossover rate is of the order of 1\nper gamete production (the highest tested).\n\nThis nonlinear relation between size of population and crossover rate could\nbe explained on the basis of the genetic structure of individual genomes in\nthe simulated populations. In Fig. 8 we have shown the frequency of\ndefective genes in the genetic pool of populations for $N_{\\max} = 10000$ under\ncrossover rates 0, 0.1 and 1. The frequency of defective genes expressed\nbefore minimum reproduction age ($R = 80$) in populations without crossover is\n0.5. Since $T = 1$, if the distribution of defects would be random the\nprobability of any individual to survive until the reproduction age $R$ would be\n$0.75^R$ (negligibly small for large $R > 30$). Thus, to survive, individuals \nhave to complete their\ngenomes of two complementing bit-strings (haplotypes). For more efficient\nreproduction the number of different haplotypes should be restricted and in\nfact there are only two different complementing haplotypes in the whole\npopulation as it was shown in \\cite{waga}. In such populations, the\nprobability of forming the offspring surviving until the reproduction age is\n0.5. Note, that recombination at any point inside the part of the genome\nexpressed before reproduction age $R$ produces a gamete which is not\ncomplementary to any other gamete produced without recombination or with\nrecombination in an other point. Thus, crossovers in such populations are\ndeleterious for the offspring.\nOn the other extreme, with crossover probability = 1, populations are under\npurifying selection. The fraction of defective genes in the population is\nkept low (about 0.1, compared with 0.5 without recombination), to enable\nthe surviving of the offspring until their reproduction period. The critical\ncrossover frequency close to 0.12 is connected with a sharp transition from\nthese two strategies of genomic evolution: complementarity and purifying\nselection. In Fig. 9 the frequency of defective genes expressed before the\nreproduction age is plotted. For lower crossover rates the fractions of\ndefective genes are kept at the level 0.5, for higher crossover rates they\nare close to 0.1. Close to the critical frequency of crossover, defective\ngenes located at both ends of the region of genomes expressed before the\nreproduction age are forced to obey the purifying selection which eliminates\nsome defects (Fig. 8).\n\nIn the case of small populations, the probability of meeting two closely related\npartners (high inbreeding coefficient) is high and as a consequence, there\nis higher probability of meeting two defective alleles in the same locus in\nzygote which determines phenotypic defect and eliminates the offspring from\nthe population. In such condition the strategy of completing the genome of\ntwo complementing haplotypes is more effective. Nevertheless, this strategy\nis not the best if effective populations are very large, with low inbreeding\ncoefficient, when the probability of meeting two identical haplotypes is\nnegligible. Thus, comparing very large populations with very small ones we\ncan observe the inbreeding depression. On the other hand, this strategy in\nsmall populations leads to the emerging of very limited number of different\nhaplotypes in the populations (in extreme only two). These haplotypes are\ncharacterized by a specific sequence of defective alleles. Independent\nsimulations generate haplotypes with different sequence of defective\nalleles. Mixing two or more populations evolving independently decreases the\nprobability of meeting in one zygote two complementing haplotypes, this\ndifference results in outbreeding depression (seen in Figs.3b and 4).\n\n\\section{Conclusion}\nWe varied the parameters of the sexual Penna ageing model, in particular by \nseparating the population into reproductively isolated groups and\/or having\nlonger bit-strings and a high minimum age of reproduction. We could observe\nand interpret inbreeding depression, outbreeding depression, and outbreeding \nadvantage, through the counterplay of purifying selection and of haplotype\ncomplementarity. Purifying selection tries to have as few mutations in \nthe bit-strings, like haplotype 00000000 for $L=8$, while haplotypes 01100101 \nand 10011010 are complementary. In both cases, deleterious effects from\nmutations are minimised.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\\subsection{Motivation}\n\\label{subsec:Motivation}\nEven though infrequent, peak load is a cause for increased capital and operating expenses for power networks. The reason for this is twofold; the higher need for grid reinforcements and the use of more expensive fossil fuel-based generators to satisfy the peak loads (for a short duration). \nThereby to manage the extra costs due to peak demand, network operators include a peak demand tariff (commensurate to the peak load) in the electricity bills of commercial customers. This motivates large-scale commercial customers such as universities and manufacturers to manage their demand better and invest in assets such as solar photovoltaic (PV) panels and stationary batteries that have the potential to shift their demand and reduce their electricity bills. A secondary effect of this is the reduction in CO\\(_2\\) emissions because of the lesser usage of fossil fuel generators, which has the potential to combat climate change and aid the process of decarbonization of our power networks.\n\nHowever, optimal scheduling of (schedulable) loads and batteries (to minimize electricity costs) requires predictions of inflexible load (baseload), solar PV generation and the spot price of electricity. The first two variables majorly depend on weather conditions. Also, increasing renewable generation in the generation mix has led to increasing price volatility in the Australian national electricity market (NEM) which also creates a correlation between electricity prices and the weather \\cite{downey2022untangling}. This makes designing such algorithms challenging as most commercial activities are planned over the mid to long term; hence require reliable predictions. Additionally, the scheduling problems are generally mixed-integer linear programs (MILPs) which are NP-hard and may be intractable depending on the formulation and problem size. \n\nSurrounding this premise, the IEEE Computational Intelligence Society (IEEE-CIS) partnered with Monash University (Victoria, Australia) to conduct a competition seeking technical solutions to manage Monash's micro-grid containing rooftop solar panels and stationary batteries~\\cite{IEEECIS}. The main challenge was to develop an optimal scheduling algorithm for Monash's lecture theatres and operation of batteries to minimize their electricity bill, considering their baseload, solar generation and NEM electricity spot prices. To this end, the contestants were provided with actual time series data of the building loads without any lecture program (baseload) and solar generation from the micro-grid. So the contestants were expected to predict the baseload and solar generation by taking into account real-world weather data (also provided) for one month in the future. Following this, the contestants had to use actual electricity spot prices for the same duration along with these predictions for the optimal scheduling algorithm. \n\n\n\n\\subsection{Related work} \n\\label{subsec:Related Work}\nWe developed separate algorithms for the solar and baseload forecasts based on practical insights and the given data. Many approaches are available in the literature for both data types, a relatively mature and crowded research field. Therefore, the methods we have developed for this problem borrow themes from this mature literature. \n\nFor solar predictions, the basic theme used is similar to the ``clear sky'' models for forecasting solar irradiance and thereby estimating PV generation \\cite{Palani2017}. The general idea here is to create a baseline model for PV generation, assuming the sky is clear, meaning there is no cloud coverage or temperature variance. This baseline is then modified based on actual or expected weather conditions to estimate the actual PV generation. The literature around this idea is based on physical models of irradiance calculations, where equations are used to develop the baseline for a given geographical location \\cite{Palani2017}. Newer methods use data-driven techniques such as time series forecasting and machine learning models~\\cite{long2014analysis}. The data-driven methods are gaining popularity because of their robustness, speed and geographic adaptability. \n\nIntuitively, this method gives reliable solar forecasts because of solar generation's seasonal and diurnal nature. However, the main drawback is the lack of data specifically related to ``clear sky\" days; Consequently, we use the most commonly occurring day's generation as baseline, which can be discovered by our previous work~\\cite{yuan2021irmac}.\n\nFor the baseload forecasting, the theme used is the application of ensemble methods, as the current state of the art identifies these methods to be most accurate when compared to stand-alone methods \\cite{raviv2015}. We use mainly a combination of random forest (RF), gradient boost (GB), autoregressive integrated moving average (ARIMA) and support vector machine (SVM), which are all well-studied and standard methods for time series forecasting \\cite{HONG2016914}. Specifically, we apply different forecasting methods to different sub-series after disaggregating the load profiles using Seasonal and Trend decomposition using Loess (STL) due to their cyclic patterns \\cite{Theodosiou2011forecasting}.\n\nThe optimal scheduling of distributed energy resources (DER) and controllable\/uncontrollable loads for micro-grids comes under the class of problems commonly referred to as energy management problems. However, these classes of problems tend to be non-convex and non-linear in nature because of the constraints associated with the scheduled devices, e.g., scheduling uninterruptible loads such as lectures that cannot be stopped once started. As reviewed by the authors in~\\cite{EMP_Review}, the majority of these problems are modeled via classical optimization approaches such as MILP~\\cite{OptimalMILP} and mixed-integer non-linear programs (MINLPs)~\\cite{MINLP_Eg}. Other approaches include dynamic programming~\\cite{DP_Eg}, rule-based optimization and meta-heuristic algorithms~\\cite{Meta_Eg}. \n\nSince these problems tend to be non-convex and non-linear, the major challenge associated with solving these problems is tractability and (in the case of meta-heuristic methods) convergence to a global optimum. The mathematical literature surrounding MILPs offers a variety of techniques to simplify these problems before solving them, and the availability of solvers such as Gurobi\\textsuperscript{\\textregistered} for MILPs make them a very attractive option for solving these problems. Also, unlike other models where the algorithms (may) converge to a local optimal solution, modern MILP solvers can obtain globally optimal solutions for this class of problems.\n\n\n\\subsection{Contributions}\n\\label{subsec:Contributions}\nBased on the related work studied, this paper offers the following contributions to the body of knowledge: \n\\begin{itemize}\n \\item A solar forecasting algorithm using training data from refined motif (RM) discovery technique and use of an over-parameterized 1D-convoluted neural network (1D-CNN) implemented via residual networks (ResNet). To overcome the lack of ``clear sky'' data, we incorporated the work done previously by Rui Yuan et al.~\\cite{yuan2021irmac} to identify RMs in the given solar generation dataset. An RM is the most repetitive pattern within a given time series, which can be extracted along with the exogenous variables (e.g., weather information) associated with this pattern. Using this as a baseline, we estimated solar generation by training an over-parameterized 1D-CNN. Some studies have shown that over-parameterization of CNNs can lead to better performance at the expense of longer training time \\cite{Balaji2021, Bubeck2021,Power2021}. Therefore, a ResNet was implemented to develop a deeper NN but with faster computation time.\n \n \\item An optimal micro-grid scheduling algorithm is solved based on real-world data for a university-based application. To the best of our knowledge, there has not yet been a study to co-optimize lectures schedule (and associated resources) and battery operation (with PV panels). Therefore, in this paper, we have developed and tested our algorithm using real-world data and practical case instances provided by Monash University. Given the large problem size (one-month schedule) and the presence of a quadratic term in the objective function for the cost of peak demand, we proposed a two sub-problem approach to formulate a tractable problem. The first sub-problem was used to limit the peak demand throughout the month, eliminating the quadratic term from the objective. Then, the peak demand was used to solve the second sub-problem to minimize the total electricity costs.\n\\end{itemize}\n\n\nThe rest of the paper is structured as follows. Section \\ref{sec:Background} introduces the data, information and problem requirements. Section \\ref{sec:Methodology} proposes the forecasting and optimization methodologies. Section \\ref{sec:Results and Discussion} presents the numerical results of the scheduling algorithm based on time series forecasting. Finally, we conclude the paper section \\ref{sec:Conclusion}.\n\n\n\n\\section{Background Information}\n\\label{sec:Background}\nThis section describes the competition requirements and data provided by the organizers.\n\n\n\\subsection{Problem statement}\n\\label{subsec:Problem Statement}\nThis section describes the scheduling constraints and objective function provided by the competition organizers and has been adapted from~\\cite{IEEECIS}. We were required to develop prediction algorithms for the baseload of six buildings in Monash University and the solar generation of PV panels connected to them. After this, we had to use these predictions to optimally schedule lecture activities and battery operation while minimizing the electricity costs (i.e., electrical energy consumption cost plus peak demand charge). We were allowed to consider the electricity prices as known parameters to simplify the problem and focus on predicting solar generation and baseload. \n\nFor each lecture activity, we are provided with several small or large rooms needed, the electrical power consumed per room and the duration of the activities (in steps of 15 minutes). We are also provided with a list of precedence activities, i.e., activities that must be performed at least one day before the activity in question. For the batteries, we are provided with the maximum energy rating, i.e., state of charge (SOC), the peak charge\/discharge power and charge\/discharge efficiency. \n\nThe scheduling must start from the first Monday of the month, and each activity must; 1) take place within office hours (9:00--17:00), 2) happen once every week and recur on the same day and time every week. Also, the number of rooms (associated with each activity) allocated at a given time interval must not exceed the total number of rooms in the six buildings. The batteries \ncan operate in one of three states; charging, discharging or idle, which is determined by the scheduling algorithm. However, the battery power must be at the peak charge\/discharge capacity once operated.\n\n\\subsection{Dataset information}\n\\label{subsec:Data set}\nThe organizers provided us with datasets for the electricity consumption of the six buildings and their connected solar panels~\\cite{IEEECIS}. About two years' worth of data was provided with a granularity of 15 minutes. However, the datasets contained many missing points, especially the demand profiles with consecutive weeks of unrecorded data, as shown in \\cite{gitlabRay}.\nInformation regarding COVID-19 lockdown and how it influenced the forecast and earlier timelines were also provided to aid in incorporating the effect of COVID-19 pandemic on baseload in our models. ERA5 climate data was also provided as exogenous variables for the time series forecasting via a partnership with OikoLabs~\\cite{IEEECIS}. \nAlso, we were allowed to use electricity price information from the Australian NEM~\\cite{NEM_Data} as a parameter for objective function calculation. The input parameters used for solar and building forecasting are summarized in Table~\\ref{tab: inputs-sensitivity}.\n\\begin{table}[!h]\n \\centering\n \\caption{Input data used for forecasting}\n \\label{tab: inputs-sensitivity}\n \\resizebox{0.5\\textwidth}{!}{\n \\begin{tabular}{|l|*{2}{c|}}\n \\hline\n \\rowcolor{AdOrange} \\backslashbox{Weather inputs (15mins)}{Forecast outputs (15mins)}\n &\\makecell{Building \\\\(15mins)}&\\makecell{Solar\\\\ (15mins)}\\\\\\hline\n \\rowcolor{AdOrangeLight}Temperature ($^\\circ$C) & \\checkmark & \\checkmark \\\\\\hline\n \\rowcolor{AdCream}Dewpoint temperature ($^\\circ$C) &\\checkmark &\\\\\\hline\n \\rowcolor{AdOrangeLight}Wind speed (m\/s) & &\\checkmark \\\\\\hline\n \\rowcolor{AdCream}Relative humidity (0-1) & \\checkmark & \\checkmark\\\\\\hline\n \\rowcolor{AdOrangeLight}Surface solar radiation (W\/$m^2$) &\\checkmark &\\checkmark\\\\\\hline\n \\rowcolor{AdCream}Total cloud cover (0-1) & &\\checkmark\\\\\\hline\n \\rowcolor{AdOrangeLight}Occupancy (0-1) & \\checkmark&\\\\\\hline\n \\rowcolor{AdCream}Annual harmonics (0-1) & &\\checkmark\\\\\\hline\n \\end{tabular}\n }\n\\end{table}\n\\subsection{Instance information}\n\\label{subsec:Instance}\n\nTo test the algorithms each team developed, the competition organisers provided ten problem instances. Each instance provided information regarding the number of activities to be scheduled and the specifics of each activity; duration of the activity in time steps, power consumed per room for each time interval, number of large\/small rooms required and a precedence list. The precedence list contains the activities that must happen one day before the specified task. Depending on the number of activities in each instance, they were divided into two sub-categories; large (200 activities) and small (50 activities) instances. We report simulation results for one large and one small instances in section~\\ref{sec:Results and Discussion}.\n\n\n\\section{Methodology}\n\\label{sec:Methodology}\nThis section describes the methodology developed to solve the problem described in~\\ref{subsec:Problem Statement}. Mainly, we outline our data cleaning strategies, forecasting methodology and scheduling algorithms used in the competition.\n\n\n\\subsection{Solar generation forecast}\n\\label{subsec:Solar generation forecast}\n\n\\begin{figure}[!ht]\n\\centerline{\\includegraphics[clip, trim=3.3cm 4.3cm 2.75cm 3.95cm,width=0.95\\linewidth]{figs\/solar_forecast_2.pdf}}\n\\caption{Block diagram of the proposed solar generation prediction model}\n\\label{fig:solar}\n\\end{figure}\n\nThe solar generation forecast methodology is depicted in Fig.~\\ref{fig:solar}. The solar data was mainly clean, with no outliers and a few missing points, which were removed during the cleaning process. Since solar time series data is repetitive in nature, we used the method described in~\\cite{yuan2021irmac} to extract the RM (most repetitive patterns) and its associated weather conditions. The intuitive idea is to train a neural network (NN) to reshape this RM with respect to the difference in weather conditions between the RM weather conditions and the actual data weather conditions. \n\n\\subsection{Baseload forecast}\n\\label{subsec:baseload forecast}\n\n\\begin{figure}[!ht]\n\\centerline{\\includegraphics[clip, trim=4.5cm 4.2cm 0.4cm 1.5cm,width=0.95\\linewidth]{figs\/baseload_forecast_1.pdf}}\n\\caption{Block diagram of the proposed baseload prediction model}\n\\label{fig:base_load}\n\\end{figure}\n\nThe block diagram of the methodology used for baseload forecasting is shown in Fig.~\\ref{fig:base_load}, which demonstrates the use of different ensemble forecasting methods for different buildings. We applied the same procedure for the data cleaning process as in the PV generation forecast. However, since there were more missing points in the demand profiles, which lasted for many weeks consecutively, we used the average value from the previous week for the long-period missing data. The choice of ensemble methods was achieved through the use of the Python sklearn package's voting regressor function \\cite{SklearnVoting}. For this dataset, this resulted in a combination of RF, SVM, ARIMA and GB methods for best accuracy. In terms of the length of training data, due to the highly dynamic nature of demand profiles, which was exacerbated during the COVID-19 period, we only used 2 months of historical data for training.\n\nDue to the seasonal\/cyclic patterns demonstrated in building 1, 2 and 5, as shown in \\cite{gitlabRay}, we utilized STL decomposition to capture these properties by disaggregating the load profiles into three components, namely trend, seasonality and residual. We then trained each profile with different methods depending on the performance of the historical data. In the end, the predictions from these three sub-series were combined to obtain the final building demand forecast. For buildings 0 and 4 with no clear cyclic patterns, we observed that either RF or GB (without STL decomposition) were sufficient to capture their repetitive nature. Lastly, for building 3, historical data showed that the load only took on a finite set of discrete values; thus, we fixed the predictions to the median.\n\n\n\\subsection{Optimal scheduling algorithm}\n\\label{subsec:Optimal scheduling}\n\n\\begin{figure}[!ht]\n\\centerline{\\includegraphics[clip, trim=3.6cm 9.3cm 1.2cm 3.5cm,width=0.95\\linewidth]{figs\/scheduling_alg_1.pdf}}\n\\caption{Block diagram of the proposed optimal scheduling algorithm}\n\\label{fig:scheduling_alg}\n\\end{figure}\n\nWe developed a MIP-based optimal scheduling problem to represent all the scheduling constraints, where the objective function was provided by the organizers as follows: \n\\begin{flalign}\n\\label{eq:QP objective}\nO(P_{t}) = \\frac{0.25}{1000}\\sum_{t \\in \\mathcal{T}}\\, P_{t} \\, \\lambda^{\\text{aemo}}_{t} + 0.005 \\, \\left(\\max_{t \\in \\mathcal{T}} P_{t}\\right)^2 &\n\\end{flalign}\n\nHere \\(P_{t}\\) and \\(\\lambda^{\\text{aemo}}_{t}\\) are the net demand and AEMO electricity spot price at time \\(t\\) over the time horizon \\(\\mathcal{T}\\) respectively. Therefore, it can be seen that the problem is a MIQP problem, which is NP-hard and hence most likely intractable given the size of this problem (\\(\\|\\mathcal{T}\\| = 2880\\)).\nHence, to improve the tractability, we reformulated the MIQP problem into two MILP sub-problems, as demonstrated in Fig.~\\ref{fig:scheduling_alg}. The first sub-problem is used to estimate an upper bound for the maximum demand (\\(P^{\\text{max}}_1 = \\max_{t \\in \\mathcal{T}} P_{t}\\)), which is used to remove the quadratic term from the objective function in the second sub-problem. \n\nThe scheduling is developed for a time duration \\(\\Delta t\\) with \\(P^{\\text{base}}_t \\text{ and } P^{\\text{solar}}_t\\) being the baseload and solar generation forecasts, respectively. For a set of batteries \\(b \\in \\mathcal{B}\\) over the time horizon \\(t \\in \\mathcal{T}\\), we define \\(E^{\\text{bat}}_{b, t}\\) as battery SOC in kWh, \\(P^{\\text{bat}}_{b, t}\\) is the battery power in kW, \\(E^{\\text{cap}}_b\\) is the battery capacity in kWh, \\(B^{\\text{c}}_{b}\/B^{\\text{d}}_{b}\\) is the battery charge and discharge power in kW, \\(u^{\\text{c}}_{b, t}\/u^{\\text{d}}_{b, t}\\) are the battery charge\/discharge status, and \\(\\eta^{\\text{c}}_{b,t}\/\\eta^{\\text{d}}_{b,t}\\) are the charging\/discharging efficiencies. For a set of activities \\(a \\in \\mathcal{A}\\) over the same time horizon, we define \\(A^{\\text{kW}}_a\\) as the power required per room for the activity in kW, \\(R^{\\text{large}}_{a}\/R^{\\text{small}}_{a}\\) are the number of large\/small rooms needed for the activity, \\(P^{\\text{sched}}_{t}\\) be the total power scheduled in kW, and \\(u^{\\text{start}}_{a,t}\/u^{\\text{active}}_{a,t}\\) are binary variables which track the activity start time and active time, respectively. \\(\\mathcal{D}_{a}\\) is the duration set for the activity, and \\(\\mathcal{P}_{a}\\) is the set of activities that must precede activity \\(a\\). Additionally, we define \\(\\mathcal{T}_{n}\\) as the set of time intervals where activities must not be active, \\(\\mathcal{T}_f\\) is the set of time intervals for the first week and \\(L^{\\text{day}}_{t}\\) corresponds to the day of the time interval \\(t\\). We present the equations used for the two sub-problems in the below sections.\n\n\\subsubsection{Sub-problem 1}\n\\label{subsubsec:Sub-problem 1}\n\\begin{subequations}\n\\begin{flalign}\n\\label{eqref:Objective 1}\n\\min_{\\psi_1} P^{\\text{max}}_{1} &\n\\end{flalign}\n\\begin{flalign}\n&\\textbf{Subject to,} \\nonumber&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Max power constraint}\n0 \\leq P_{t}\\leq P^{\\text{max}}_{1} \\quad \\forall \\, t \\in \\mathcal{T}& \n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:SOC limit}\n0 \\leq E^{\\text{bat}}_{b, t} \\leq E^{\\text{cap}}_{b} \\quad \\forall \\, t \\in \\mathcal{T}, \\, b \\in \\mathcal{B}&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Charge-Discharge}\nu^{\\text{c}}_{b, t} + u^{\\text{d}}_{b, t} \\leq 1 \\quad \\forall \\, t \\in \\mathcal{T}, \\, b \\in \\mathcal{B}&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Battery power}\nP^{\\text{bat}}_{b, t} = u^{\\text{c}}_{b, t}\\,\\frac{B^{\\text{c}}_{b}}{\\sqrt{\\eta_{b}^{\\text{c}}}} - u^{\\text{d}}_{b, t}\\,B^{\\text{d}}_{b}\\,\\sqrt{\\eta_{b}^{\\text{d}}} \\quad \\forall \\, t \\in \\mathcal{T}, \\, b \\in \\mathcal{B} &\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Battery capacity}\nE^{\\text{bat}}_{b, -1} = E^{\\text{cap}}_{b} \\quad \\forall \\,b \\in \\mathcal{B}&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:SOC value}\nE^{\\text{bat}}_{b, t} = E^{\\text{bat}}_{b, t-1} + \\Bigg(u^{\\text{c}}_{b, t}\\, B^{\\text{c}}_{b} - u^{\\text{d}}_{b, t}\\, B^{\\text{d}}_{b} \\Bigg) \\Delta t \\quad \\forall \\, t \\in \\mathcal{T}, \\, b \\in \\mathcal{B}& \n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Scheduled power}\nP^{\\text{sched}}_{t} = \\sum_{a \\in \\mathcal{A}} u^{\\text{active}}_{a,t} \\, A^{kW}_a\\, \\Bigg(R^{\\text{small}}_{a} + R^{\\text{large}}_{a}\\Bigg) \\quad \\forall \\, t \\in \\mathcal{T}&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Net export}\nP_{t} = P^{\\text{base}}_{t} - P^{\\text{solar}}_t + P^{\\text{sched}}_{t}+ \\sum_{b \\in \\mathcal{B}}P^{\\text{bat}}_b \\quad \\forall \\, t \\in \\mathcal{T} &\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Room requirement}\n\\sum_{a \\in \\mathcal{A}} u^{\\text{active}}_{a, t}\\, R^{\\text{large}}_{a} \\leq H^{\\text{large}}; \\sum_{a \\in \\mathcal{A}} u^{\\text{active}}_{a, t}\\, R^{\\text{small}}_{a} \\leq H^{\\text{small}} \\; \\forall t \\in \\mathcal{T}&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Non active time}\nu^{\\text{active}}_{a, t} = 0 ; u^{\\text{start}}_{a, t} = 0 \\quad \\forall \\, t \\in \\mathcal{T}_{n}, \\, a \\in \\mathcal{A}& \n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Happen once a week}\n\\sum_{t \\in \\mathcal{T}_f} u^{\\text{start}}_{a, t} = 1 \\quad \\forall \\, a \\in \\mathcal{A} &\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Duration constraint A}\nu^{\\text{start}}_{a, t} \\leq u^{\\text{active}}_{a, t + k} \\: \\forall \\, & t \\in \\{0, 1,\\dots, \\left|\\mathcal{T}\\right| - \\left|\\mathcal{D}_{a}\\right| - 1\\}, \\, k \\in \\mathcal{D}_{a}, a \\in \\mathcal{A}&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Duration constraint B}\n\\sum_{t \\in \\mathcal{T}_f} u^{\\text{active}}_{a, t} = \\left| \\mathcal{D}_{a} \\right| \\quad \\forall \\, a \\in \\mathcal{A}&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Precedence constraints}\n\\sum_{t \\in \\mathcal{T}_f} \\left(u^{\\text{active}}_{a, t}\\, L^{\\text{day}}_{t} - u^{\\text{active}}_{k, t}\\, L^{\\text{day}}_{t}\\right) \\geq 1 \\; \\; \\forall \\, k \\in \\mathcal{P}_{a}, a \\in \\mathcal{A}&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Start at same time every week}\nu^{\\text{start}}_{a, t} = u^{\\text{start}}_{a, t - \\alpha} \\quad \\forall \\, t \\in \\mathcal{T}_{f}^{c}, \\, a \\in \\mathcal{A} &\n\\end{flalign}\n\\end{subequations}\n\n\\noindent where \\(\\psi_1 = \\{P, P^{\\text{bat}}, P_1^{\\text{max}}, E^{\\text{bat}},u^{\\text{c}}, u^{\\text{d}}, u^{\\text{active}}, u^{\\text{start}} \\}\\). Equation~\\eqref{eqref:Max power constraint} helps find the peak demand of the schedule and prevents energy export to the grid. Equations~\\eqref{eqref:SOC limit}--\\eqref{eqref:SOC value} represent battery related constraints. Equations~\\eqref{eqref:Scheduled power}--\\eqref{eqref:Net export} are used to calculate the net energy imported from the grid, and equation~\\eqref{eqref:Room requirement} ensures the number of rooms engaged through activities is within the limit. Equation~\\eqref{eqref:Non active time} ensures activities do not start outside office hours and before the first Monday of the month, equation~\\eqref{eqref:Happen once a week} ensures activities happen once a week, equations~\\eqref{eqref:Duration constraint A}--\\eqref{eqref:Duration constraint B} ensure activities are active for the necessary duration after starting, equation~\\eqref{eqref:Precedence constraints} satisfies that all precedence (prerequisite) activities happen a day before and equation~\\eqref{eqref:Start at same time every week} ensures activities recur at the same time every week. The objective here is to minimize peak demand while satisfying all the scheduling constraints. \n\nSince this problem is only used to solve for the upper bound maximum demand \\(P^{\\text{max}}_1\\), it is not necessary to solve the full problem. We use the property that the solution of the LP relaxation of a minimization-based MILP problem is its lower bound (\\(y^{*}_{\\text{MILP}} \\geq y^{*}_{\\text{LP}}\\))~\\cite{Vielma2015} to just solve the LP relaxation and use this for sub-problem 2, thereby improving tractability of our method.\n\n\n\\subsubsection{Sub-problem 2}\n\\label{subsubsec:Sub-problem 2}\n\\begin{subequations}\n\\begin{flalign}\n\\min_{\\psi_2} \\frac{0.25}{1000}\\sum_{t \\in \\mathcal{T}}\\, P_{t} \\, \\lambda^{\\text{aemo}}_{t} &\n\\end{flalign}\n\\begin{flalign}\n&\\textbf{Subject to, Equations~}\\eqref{eqref:SOC limit} \\text{\\---} \\eqref{eqref:Precedence constraints}\\nonumber&\n\\end{flalign}\n\\begin{flalign}\n\\label{eqref:Load limit}\n0 \\leq P_{t} \\leq \\beta^{\\text{mul}} \\,P^{\\text{max}^{*}}_{1} &\n\\end{flalign}\n\\end{subequations}\n\n\\noindent where \\(\\psi_2 = \\left(\\psi_1 \\setminus \\{P^{\\text{max}}_{1}\\} \\right)\\). Using equation~\\eqref{eqref:Load limit}, we ensure that peak demand does not exceed \\(\\beta^{\\text{mul}} P^{\\text{max}^{*}}_1\\) limiting the demand charge. Intuitively, \\(\\beta^{\\text{mul}} \\geq 1\\) is a multiplier used to change the tightness of the MILP, creating a trade-off between computation time and electricity cost; higher values of \\(\\beta^{\\text{mul}}\\) decreases computation speed but increases the electricity cost (objective of the problem) and vice versa. \n\n\\begin{figure*}[t]\n \\centering\n \\subfigure[Difference in the RM at and predicted generation weather conditions for Nov \\(4^{\\text{th}}\\) 2020]\n {\n \\label{subfig:RM differences}\n \\includegraphics[width=0.29\\textwidth]{figs\/RM_Exogenous.pdf} \n } \n %\n \\subfigure[RM vs. predicted solar generation for Nov \\(4^{\\text{th}}\\) 2020]\n {\n \\label{subfig:RM output}\n \\includegraphics[width=0.29\\textwidth]{figs\/RM_Output.pdf} \n } \n %\n \\subfigure[Scheduled power, predicted net load and NEM spot price for Nov \\(2^{\\text{nd}}\\) 2020]\n {\n \\label{subfig:Optimization output}\n \\includegraphics[width=0.33\\textwidth]{figs\/Competition_Results.pdf}\n }\n \\vspace{-0.4em}\n \\caption{Prediction and optimization results}\n \\label{fig:Figure ref}\n\\end{figure*}\n\n\\section{Results and Discussion}\n\\label{sec:Results and Discussion}\n\nThe simulations were run for November 2020 as per the competition stipulation with \\(\\Delta t = 15 \\text{ min}, \\\n\\|\\mathcal{T}\\|=2880\\) and other instance information as specified in \\ref{subsec:Instance}. For the scheduling, two \\(\\beta^{\\text{mul}}\\) values were chosen; 1.10 for small instances and 1.15 for large instances. These values were obtained via trial and error. The prediction and optimization algorithms were run in Python 3.8.8 using Gurobi\\textsuperscript{\\textregistered} solver. \n\n\\begin{table}[t]\n \\caption{Baseload and PV generation prediction NRMSE}\n \\centering\n \\begin{tabular}{ccc|ccc}\n \\toprule\n \\rowcolor{AdOrange} \\textbf{Building \\#} & \\textbf{Proposed} &\n \\textbf{BI} & \\textbf{Solar \\#} & \\textbf{Proposed} & \\textbf{BI} \\\\\n \\midrule\n 0 & \\textbf{0.43} & 0.43 & 0 & \\textbf{0.72} & 0.81 \\\\\n \\rowcolor{AdOrangeLight} 1 & \\textbf{0.31} & 0.34 & 1 & \\textbf{0.77} & 0.77 \\\\\n 2 & \\textbf{0.23} & 0.24 & 2 & \\textbf{0.74} & 0.79 \\\\\n \\rowcolor{AdOrangeLight} 3 & \\textbf{0.32} & 0.32 & 3 & \\textbf{0.74} & 0.75 \\\\\n 4 & \\textbf{0.40} & 0.40 & 4 & \\textbf{0.78} & 0.79 \\\\\n \\rowcolor{AdOrangeLight} 5 & \\textbf{0.13} & 0.14 & 5 & \\textbf{0.72} & 0.75 \\\\\n \\bottomrule\n \\end{tabular}\n \\label{tab:mae}\n\\end{table}\n\nSince each of the twelve datasets has different scales, we decide to use the normalised root mean square error (NRMSE), which has been used to facilitate the comparison among time series with varying magnitudes \\cite{das2018evaluation}. The metric is calculated as follows:\n\\begin{flalign}\n\\label{eq:NRMSE}\n& \\text{NRMSE} = \\sqrt{\\sum_t \\frac{\\hat{y}_t - y_t}{\\|\\mathcal{T}\\|}} \\times \\frac{1}{\\bar{y}}\n\\end{flalign}\nwhere $\\hat{y}_t$ is the model predictions, $y_t$ is the model actual data and $\\bar{y}$ represents the mean of the actual data.\n\nIn order to demonstrate the superiority of our methods, we compare our proposed techniques with the best individual (BI) prediction models that are mentioned in section \\ref{subsec:Related Work}. Please note that the BI models are picked from out-of-sample validation on the actual datasets using the same inputs and training length as the proposed methods. Table~\\ref{tab:mae} displays the NRMSE of the actual data in November 2020 between the proposed and the BI models. It can be seen that the proposed models achieve lower errors in almost all the twelve profiles except for the demand of buildings 0 and 4. As mentioned in section \\ref{subsec:baseload forecast}, these two profiles are predicted using either RF or GB due to their non-cyclic behavior. Therefore, they have the highest errors among all the demand profiles. For the PV generation time series, the proposed method also achieves better prediction accuracy with quite consistent NRMSE due to the high dependence of solar PV generation on weather data.\n\nFig.~\\ref{subfig:RM differences} and \\ref{subfig:RM output} are used to explain the intuition behind the RM method for the PV generation; the former shows the difference between RM \nand a predicted generation weather condition for November 4\\textsuperscript{th}, 2020 while the latter shows the RM and the predicted generation profiles for the same day. We can see that as the difference in temperature is positive, the predicted generation is lower than the RM and vice versa. The ResNet is then trained with the changes in weather conditions to learn how to reshape the RM to predict solar generation.\n\nFig.~\\ref{subfig:Optimization output} shows the power scheduled for the small and large instances, the predicted load and the NEM spot price profiles for Novemeber \\(2^{\\text{nd}}\\) 2020. It can be seen that during lower price periods (intervals 24--72), most of the load is scheduled and the scheduled load is clipped to the value obtained from sub-problem 1. During high price periods (intervals 72--88), the scheduled load is lower than the predicted baseload due to battery operation. This ensures that electricity cost is minimized. \n\n\\begin{table}[b]\n \\caption{Electricity cost from different forecasting results}\n \\centering\n \\begin{tabular}{c|ccc}\n \\toprule\n \\rowcolor{AdOrange} \\textbf{Forecast methods} &\n \\textbf{Base case} & \\textbf{BI} & \\textbf{Proposed} \\\\\n \\midrule\n \\rowcolor{AdOrangeLight} \\textbf{Cost (\\$AUD)} & 589356 & 466652 & \\textbf{357210} \\\\\n \\bottomrule\n \\end{tabular}\n \\label{tab:opti_cost}\n\\end{table}\n\nThe organizers provided a base schedule\/electricity cost for the problem based on no predictions and an \\(\\epsilon\\)-greedy algorithm \\cite{IEEECIS}. Table~\\ref{tab:opti_cost} compares the electricity cost from three different forecast methods, where the base case represents the electricity cost provided by the organizers; whereas BI and the proposed cost represent the cost optimized on their respective forecasts. Our methodology reduced this cost by roughly 23\\% and 40\\% compared to the BI and base case, respectively. We were able to solve small instances (158400 binaries) in an average time of 447 seconds per instance and large instances (590400 binaries) in 1550 seconds per instance. It can be seen that even though the problem size grew 3.73 times, the time to solve the problem increased almost linearly and not exponentially with problem size as is the case in MIPs. It further demonstrates the tractability and scalability of our algorithm. \n\n\n\\section{Conclusion}\n\\label{sec:Conclusion}\n\nIn this paper, we developed an optimal scheduling solution for solving a real-world electricity cost minimization problem for six buildings at Monash University in Melbourne, Australia. A two-step solution was presented in this paper, which included forecasting and an optimization step. Our solution with energy forecast reduced the total cost of electricity roughly by 40\\%. We would like to highlight that we placed fifth in the competition among 50 participants and were off from the best solution only by 4\\% with respect to the objective value due to their better forecast.\nFor future work, the forecasting of baseload and solar generation can be improved by using\nmodern machine learning techniques, such as transformers and info-GANs.\nFor the optimization, a major improvement can be devising an algorithm to select optimal \\(\\beta^{\\text{mul}}\\) considering the trade-off between the tightness (complexity) of the problem and cost minimization. \n\\bibliographystyle{IEEEtran}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\n\\section{Introduction}\n Cooperative communications via relaying is an effective measure to\n ensure reliability, provide high throughput and extend network\n coverage. It has been intensively studied in LTE-Advanced \\cite{4G-relay} and will continue to play an important role in\n the future fifth generation wireless networks. In the conventional\n two-hop one-way relay channel (OWRC), the communication takes two orthogonal phases, and suffers from the loss of spectral efficiency because of the\n inherent half-duplex (HD) constraint at the relay. Two-way relaying using the principle of physical layer network coding is proposed to\n recover the loss in OWRC by allowing the two sources to exchange\n information more efficiently \\cite{KRI}\\cite{Zhang-2Phase}\\cite{Qiang-Li-2009}. In the first phase of the two-way relay channel (TWRC), both sources\n transmit signals to the relay. In the second phase, the relay does not separate signals\n but rather broadcasts the processed mixed signals to both sources.\n Each source can subtract its own message from the received signal\n then decodes the information from the other source. The benefit of the TWRC is that,\n using the same two communication phases as the OWRC, bi-directional communication is achieved.\n\n However, note that the relay in the TWRC still operates in the HD mode thus two communication phases\n are needed. Motivated by this observation and thanks to the emerging full-duplex (FD) techniques,\n we aim to study the potential of the FD operation\n in the TWRC to enable simultaneous information\n exchange between the two sources, i.e., only one communication phase is required.\n In the proposed scheme, all nodes, including the two sources\nand the relay, work in the FD mode so they can transmit and receive\nsignals simultaneously \\cite{BLISS}\\cite{BLISS1}. However, the major\nchallenge in the FD operation is that the power of the\nself-interference (SI) from the relay output could be over 100 dB\nhigher than that of the signal received from distance sources and\nwell exceeds the dynamic range of the analog-to-digital converter\n\\cite{RII3,RII2,DAY}. {\\bl Therefore it is important that the SI is\nsufficiently mitigated. The SI cancellation can be broadly\ncategorized as passive cancellation and active cancellation\n\\cite{Sabharwal-inband}. Passive suppression is to isolate the\ntransmit and receive antennas using techniques such as directional\nantennas, absorptive shielding and cross-polarization\n\\cite{Sabharwal-passive}. Active suppression is to exploit a node's\nknowledge of its own transmit signal to cancel the SI, which\ntypically includes analog cancellation, digital cancellation and\nspatial cancellation. Experimental results are reported in\n\\cite{Duarte-Exp} that the SI can be cancelled to make FD wireless\ncommunication feasible in many cases. } In the recent work\n\\cite{Stanford-FD} and \\cite{Stanford-FD-MIMO}, the promising\nresults show that the SI can be suppressed to the noise level in\nboth single-antenna and multi-antenna cases. Below we will provide a\nreview on the application of the FD relaying to both the OWRC and\nTWRC.\n\n\n\\subsection{Literature on the FD relaying}\n\n\\subsubsection{OWRC}\nThe FD operation has attracted lots of {\\bl research interest} for\nrelay assisted cooperative communication. It is shown in \\cite{RII3}\nthat the FD relaying is feasible even in the presence of the SI and\ncan offer higher capacity than the HD mode. Multiple-input\nmultiple-output (MIMO) techniques provide an effective means to\nsuppress the SI in the spatial domain\n\\cite{RII1}-\\cite{MIMO-relay-HD}. The authors of \\cite{RII1} analyze\na wide range of SI mitigation measures when the relay has multiple\nantennas, including natural isolation, time-domain cancellation and\nspatial domain suppression. { The FD relay selection is studied in\n\\cite{KRIKIDIS} for the amplify-and-forward (AF) cooperative\nnetworks.} With multiple transmit or receive antennas at the\nfull-duplex relay, precoding at the transmitter and decoding at the\nreceiver can be jointly optimized to mitigate the SI effects. The\njoint precoding and decoding design for the FD relaying is studied\nin \\cite{RII1}, where both the zero forcing (ZF) solutions and the\nminimum mean square error (MMSE) solutions are discussed. When\nonly the relay has multiple antennas, a joint design of ZF precoding\nand decoding vectors is proposed in \\cite{CHOI_EL} to null out the\nSI at the relay. However, this design does not take into account the\nend-to-end (e2e) performance. A gradient projection method is\nproposed in \\cite{Park-null-12}, to optimize the precoding and\ndecoding vectors considering the e2e performance. When all terminals\nhave multiple antennas, the e2e performance is optimized in\n\\cite{MIMO-relay-HD} where the closed-form solution for\nprecoding\/decoding vectors design as well as diversity analysis are\nprovided.\n\n\\subsubsection{TWRC}\n\n In the early work of TWRC, the FD operation is often employed to\n investigate the capacity region from the information-theoretic viewpoint without considering the effects of the SI \\cite{capacity-twrc}\\cite{capacity-twrc-half-bit}.\n Only recently, the SI has been taken into account in the FD TWRC.\n In \\cite{FD-two-phase}, the FD operation is introduced to the relay but two one-way relaying phases are required to achieve the two-way communication to avoid\n interference.\n A physical layer network coding FD TWRC is proposed in \\cite{FD-practical}, where bit error\nrate is derived. The optimal relay selection scheme is proposed and\nanalyzed in \\cite{Song-FD-TW-Selection} for the FD TWRC using the AF\nprotocol. Transmit power optimization among the source nodes and the\nrelay node is studied in \\cite{FD-TW-SISO}, again using the AF\nprotocol. However, all existing works are restricted to the single\ntransmit\/receive antenna case at the relay thus the potential of\nusing multiple antennas to suppress the SI and improve the e2e\nperformance for the TWRC is not fully explored yet.\n\n\\subsection{Our work and contribution}\n In this work, we study the potential of the MIMO FD operation\n in the TWRC where the relay has multiple transmit\/receive antennas\n and employs the AF protocol and the principle of physical layer network coding. The two sources have single transmit\/receive antenna.\n We jointly optimize the relay beamforming matrix and the power\n allocation at the sources to maximize the e2e performance.\n To be specific, we study two problems: one is to find the achievable rate\n region and the other is to maximize the sum rate of the FD TWRC.\n Our contributions are summarized as follows:\n \\begin{itemize}\n \\item We derive the signal model for the MIMO FD TWRC and\n propose to use the ZF constraint at the relay to greatly simplify the model\n and the problem formulations.\n \\item We find the rate region by maximizing one source's rate\n with the constraint on the other source' minimum rate. We propose an\n iterative algorithm together with 1-D search to find the local\n optimum solution. At each iteration, we give analytical solutions for the transmit beamforming vector and the power control.\n \\item We tackle the sum rate maximization problem by employing a similar iterative algorithm.\n At each iteration, we propose to use the DC (difference of convex functions) approach to optimize the transmit beamforming vector and we solve the power allocation analytically.\n \\item We conduct intensive simulations to compare the proposed FD\n scheme with three benchmark schemes and clearly show its\n advantages of enlarging the rate region and improving the sum rate.\n \\end{itemize}\n\n\n\n\nThe rest of the paper is organized as follows. In Section II, we\npresent the system model, the explicit signal model and problem\nformulations. In Section III, we deal with the problem of finding\nthe achievable rate region; in Section IV, we address the problem of\nmaximizing the sum rate. Three benchmark schemes are introduced in\nSection V. Simulation results are presented in Section VI. Section\nVII concludes this paper and gives future directions.\n\n\n\n\n \\noindent \\emph{Notation}: The lowercase and\nuppercase boldface letters (e.g., ${\\bf x}$ and ${\\bf X}$) indicate column\nvectors and matrices, respectively. ${\\bf X}\\in \\mathbb{C}^{M\\times N}$\nmeans a complex matrix ${\\bf X}$ with dimension of $M\\times N$.\n ${\\bf I}$ is the identity matrix.\n We use $(\\cdot)^\\dagger$ to denote the conjugate\ntranspose, $\\textsf{trace}(\\cdot)$ is the trace operation, and $\\|\\cdot\\|$ is\nthe Frobenius norm. $|\\cdot|$ represents the absolution value of a\nscalar. ${\\bf X}\\succeq {\\bf 0}$ denotes that the Hermitian matrix ${\\bf X}$\nis positive semidefinite. The expectation operator is denoted by\n$\\mathcal{E}(\\cdot)$. Define $\\Pi_{\\bf X} =\n{\\bf X}({\\bf X}^\\dag{\\bf X})^{-1}{\\bf X}^\\dag$ as the orthogonal projection onto the\ncolumn space of ${\\bf X}$; and $\\Pi_{\\bf X}^\\bot = {\\bf I} - \\Pi_{\\bf X}$ as the\northogonal projection onto the orthogonal complement of the column\nspace of ${\\bf X}$.\n\n\n\n\n\n\\section{System model, signal model and problems statement}\n\\subsection{System model}\\label{GENERAL_system_model}\nConsider a three-node MIMO relay network consisting of two sources\n${\\tt A}$ and ${\\tt B}$ who want to exchange information with the aid of a\nMIMO relay ${\\tt R}$, as depicted in Fig. 1. {\\bl There is no direct\nlink between the two sources because of deep fading or heavy\nshadowing, so their communication must rely on ${\\tt R}$.} All nodes\nwork in the FD mode. To enable the FD operation, each source is\nequipped with two groups of RF chains and corresponding antennas,\ni.e., one for transmitting and one for receiving signals\\footnote{It\nis also possible to realize the FD operation using a single antenna,\nsee \\cite{Stanford-FD}.}. We assume that each source has one\ntransmit antenna and one receive antenna. We use $M_{T}$ and $M_{R}$\nto denote the number of transmit and receive antennas at ${\\tt R}$,\nrespectively. We assume $M_T>1$ or $M_R>1$ to help suppress the\nresidual SI at ${\\tt R}$ in the spatial domain.\n\\begin{figure}[t]\n\\centering\n\\includegraphics[width=0.9\\linewidth]{MIMO_FD_Relay}\n \\caption{Full-duplex TWRC with two sources and a MIMO\nrelay. The dashed line denotes the residual self-interference.}\n\\label{FIGG1}\n\\end{figure}\n We use ${\\bf h}_{XR}\\in \\mathbb{C}^{M_R\\times 1}$ and\n${\\bf h}_{RX}\\in\\mathbb{C}^{M_T\\times 1}$ to denote the directional\nchannel vectors between the source node ${\\tt X}$'s (${\\tt X}\\in\n\\{{\\tt A},{\\tt B}\\}$) transmit antenna to ${\\tt R}$'s receive antenna(s), and\nbetween ${\\tt R}$'s transmit antenna(s) to ${\\tt X}$'s receive antenna,\nrespectively. In general, channel reciprocity does not hold, e.g.,\n${\\bf h}_{XR}\\ne {\\bf h}_{RX}$, due to different transmit and receive\nantennas used. In addition, $h_{AA}, h_{BB}$ and\n${\\bf H}_{RR}\\in\\mathbb{C}^{M_R\\times M_T}$ denote the residual SI\nchannels after the SI cancellation scheme is applied\n at the corresponding nodes \\cite{KRIKIDIS,RII3}. {\\bl The statistics of\nthe residual SI channel are not well understood yet \\cite{SI-stat}.\nThe experimental-based study in \\cite{Duarte-thesis} has\ndemonstrated that the amount of SI suppression achieved by a\ncombined analog\/digital cancellation technique is a complicated\nfunction of several system parameters. Therefore for simplicity, in\nthis paper, we model each element of the residual SI channel as a\nGaussian distributed random variable} with zero mean and variance\n$\\sigma^2_X, {\\tt X}\\in \\{{\\tt A},{\\tt B},{\\tt R}\\}$. All channel links between two\nnodes are subject to independent flat fading. We assume that the\nchannel coefficients between different nodes remain constant within\na normalized time slot, but change independently from one slot to\nanother according to a complex Gaussian distribution with zero mean\nand unit variance. The global channel state information (CSI) is\navailable at the relay where the optimization will be performed. As\nwill be seen in Fig. \\ref{fig:sumrate:csi}, this is critical for the\nrelay to adapt its beamforming matrix and for the two sources to\nadjust their transmit power. The noise at the node ${\\tt X}$'s (${\\tt X}\\in\n\\{{\\tt A},{\\tt B}, {\\tt R}\\}$) receive antenna(s) is denoted as $n_X$ (${\\bf n}_X$)\nand modeled as complex additive white Gaussian noise with zero mean\nand unit variance.\n\nSince each source has only a single transmit and receive antenna, a\nsingle data stream $S_A$ and $S_B$ are transmitted from ${\\tt A}$ and\n${\\tt B}$, with the average power $p_A$ and $p_B$, respectively. To keep\nthe complexity low, ${\\tt R}$ employs linear processing, i.e., the AF\nprotocol with an amplification matrix ${\\bf W}$, to process the\nreceived signal. The node ${\\tt X}$ has maximum transmit power\nconstraint $P_X$. We will jointly optimize the transmit power of the\ntwo source nodes, $p_A$ and $p_B$ together with the amplification\nmatrix at the relay to maximize the e2e system performance.\n\n\n {\\bbl The overhead due to the CSI acquisition at each node is analyzed as follows. The received CSI\n${\\bf h}_{AR}$ and ${\\bf h}_{BR}$ at ${\\tt R}$ can be estimated by ${\\tt A}$ and\n${\\tt B}$ each sending one pilot symbol separately. The SI channel\n${\\bf h}_{RR}$ can be estimated at ${\\tt R}$ by itself sending an\n$M_T$-symbol pilot sequence. The SI channel $h_{AA}$ and $h_{BB}$\n are estimated similarly and then sent back to\n${\\tt R}$. Regarding the transmit CSI ${\\bf h}_{RA}$ and ${\\bf h}_{RB}$, ${\\tt R}$\nfirst broadcasts an $M_T$-symbol pilot sequence (this can be used\nfor the estimation of ${\\bf h}_{RR}$ simultaneously), then ${\\tt A}$ and\n${\\tt B}$ feedback their estimation to ${\\tt R}$. In addition, after ${\\tt R}$\nperforms the optimization and obtains ${\\bf W}$, $p_A$ and $p_B$, it\ntransmits ${\\bf h}_{RA}^\\dag{\\bf W} {\\bf h}_{AR}$ and ${\\bf h}_{RB}^\\dag{\\bf W}\n{\\bf h}_{BR}$ to ${\\tt A}$ and ${\\tt B}$, respectively, such that they can\ncancel their previously sent symbols. Finally, ${\\tt R}$ informs ${\\tt A}$\nand ${\\tt B}$ their transmit power $p_A$ and $p_B$, respectively.}\n\n\n\n\\subsection{Signal model}\n We assume that the processing delay at ${\\tt R}$ is given by {\\bbl a\n $\\tau$-symbol duration}, which refers to the required\nprocessing time in order to implement the FD operation \\cite{RII1}.\n$\\tau$ typically takes integer values. The delay is short enough\ncompared to a time slot which has a large number of data symbols,\ntherefore its effect on the achievable rate is negligible. At the\ntime ({\\bbl symbol}) instance $n$, the received signal ${\\bf r}[n]$ and\nthe transmit signal ${\\bf x}_R[n]$ at ${\\tt R}$ can be written as\n \\begin{equation}} \\newcommand{\\ee}{\\end{equation}\\label{eqn:rn}\n {\\bf r}[n] = {\\bf h}_{AR}S_A[n] + {\\bf h}_{BR}S_B[n] + {\\bf H}_{RR} {\\bf x}_R[n] + {\\bf n}_R[n],\n \\ee \\vspace{-3mm} and\n \\begin{equation}} \\newcommand{\\ee}{\\end{equation}\\label{eqn:xn}\n {\\bf x}_R[n] = {\\bf W}{\\bf r}[n-\\tau],\n \\ee\nrespectively. Using (\\ref{eqn:rn}) and (\\ref{eqn:xn}) recursively,\nthe overall relay output can be rewritten as \\small\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:xn2}\n {\\bf x}_R[n] &=& {\\bf W} {\\bf h}_{AR}S_A[n-\\tau] + {\\bf W}{\\bf h}_{BR}S_B[n-\\tau] \\notag\\\\ &&+ {\\bf W}{\\bf H}_{RR} {\\bf x}_R[n-\\tau] + {\\bf W}{\\bf n}_R[n-\\tau]\\notag\\\\\n &=& {\\bf W}\\sum_{j=0}^{\\infty}\\left({\\bf H}_{RR}{\\bf W}\\right)^j \\Big({\\bf h}_{AR}S_A[n-j\\tau-\\tau] \\notag\\\\ && + {\\bf h}_{BR}S_B[n-j\\tau-\\tau]+ {\\bf n}_R[n-j\\tau-\\tau] \\Big),\n \\eea \\normalsize\n {\\bl where $j$ denotes the index of the delayed symbols}.\n\n Its covariance matrix is given by\\small\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n &&{\\mathcal{E}}[{\\bf x}_R{\\bf x}_R^\\dag] = P_A {\\bf W} \\sum_{j=0}^{\\infty}\\left({\\bf H}_{RR}{\\bf W}\\right)^j {\\bf h}_{AR}{\\bf h}_{AR}^\\dag\n \\left( \\left({\\bf H}_{RR}{\\bf W}\\right)^j\\right)^\\dag{\\bf W}^\\dag \\notag\\\\\n &&+P_B {\\bf W} \\sum_{j=0}^{\\infty}\\left({\\bf H}_{RR}{\\bf W}\\right)^j {\\bf h}_{BR}{\\bf h}_{BR}^\\dag\n \\left( \\left({\\bf H}_{RR}{\\bf W}\\right)^j\\right)^\\dag{\\bf W}^\\dag\\notag\\\\\n && + {\\bf W} \\sum_{j=0}^{\\infty} ({\\bf H}_{RR}{\\bf W}\\qW^\\dag {\\bf H}_{RR}^\\dag )^j {\\bf W}^\\dag\\notag\\\\\n &&= P_A {\\bf W} \\sum_{j=0}^{\\infty}\\left({\\bf H}_{RR}{\\bf W}\\right)^j {\\bf h}_{AR}{\\bf h}_{AR}^\\dag\n \\left( \\left({\\bf H}_{RR}{\\bf W}\\right)^j\\right)^\\dag{\\bf W}^\\dag \\notag\\\\\n &&+ P_B {\\bf W} \\sum_{j=0}^{\\infty}\\left({\\bf H}_{RR}{\\bf W}\\right)^j {\\bf h}_{BR}{\\bf h}_{BR}^\\dag\n \\left( \\left({\\bf H}_{RR}{\\bf W}\\right)^j\\right)^\\dag{\\bf W}^\\dag\\notag\\\\\n &&+{\\bf W}({\\bf I}-{\\bf H}_{RR}{\\bf W}\\qW^\\dag{\\bf H}_{RR}^\\dag)^{-1}{\\bf W}^\\dag.\n\\eea\\normalsize\n\n Note that ${\\tt R}$'s transmit signal covariance $\\mathcal{E}[{\\bf x}_R{\\bf x}_R^\\dag]$, and in turn the transmit power and the SI power, are complicated functions of ${\\bf W}$, which makes the optimization problems difficult.\n To simplify the signal model and make the optimization problems more tractable, we add the ZF constraint such that the optimization of ${\\bf W}$\n nulls out the residual SI from the relay output to relay\ninput. To realize this, it is easy to check from (\\ref{eqn:xn2})\nthat the following condition is sufficient, \\begin{equation}} \\newcommand{\\ee}{\\end{equation}\n {\\bf W}{\\bf H}_{RR}{\\bf W}={\\bf 0}.\n \\ee\n Consequently, (\\ref{eqn:xn2}) becomes\n {\\small\\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:xn3}\n {\\bf x}_R[n] = {\\bf W} \\left({\\bf h}_{AR}S_A[n-\\tau] + {\\bf h}_{BR}S_B[n-\\tau] +\n {\\bf n}_R[n-\\tau]\\right),\n \\eea}\n with the covariance matrix\n {\\small\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n {\\mathcal{E}}[{\\bf x}_R{\\bf x}_R^\\dag] = p_A {\\bf W}{\\bf h}_{AR}{\\bf h}_{AR}^\\dag{\\bf W}^\\dag + p_B {\\bf W}{\\bf h}_{BR}{\\bf h}_{BR}^\\dag{\\bf W}^\\dag + {\\bf W}\\qW^\\dag.\n\\eea}\n The relay output power is\n \\begin{align}\n p_R &= \\textsf{trace}( {\\mathcal{E}}[{\\bf x}_R{\\bf x}_R^\\dag])\\notag\\\\\n &= p_A\\|{\\bf W} {\\bf h}_{AR}\\|^2 + p_B\\|{\\bf W} {\\bf h}_{BR}\\|^2 + \\textsf{trace}({\\bf W}\\qW^\\dag).\n \\end{align}\n\nThe received signal at the source ${\\tt A}$ can be written as\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n r_A[n] &=& {\\bf h}_{RA}^\\dag{\\bf x}_R[n] + h_{AA}S_A[n] + n_A[n]\\notag\\\\\n &=& {\\bf h}_{RA}^\\dag{\\bf W} {\\bf h}_{AR}S_A[n-\\tau] +{\\bf h}_{RA}^\\dag{\\bf W} {\\bf h}_{BR}S_B[n-\\tau] \\notag\\\\&& + {\\bf h}_{RA}^\\dag{\\bf W}{\\bf n}_R[n] + h_{AA}S_A[n] + n_A[n].\n \\eea\nAfter cancelling its own transmitted signal $S_A[n-\\tau]$, it\nbecomes\\footnote{Note that different from $S_A[n-\\tau]$, $S_A[n]$\ncan not be completely cancelled due to the simultaneous\ntransmission, which is the main challenge of the FD radio.}\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n r_A[n] &=& {\\bf h}_{RA}^\\dag{\\bf W} {\\bf h}_{BR}S_B[n-\\tau] + {\\bf h}_{RA}^\\dag{\\bf W}{\\bf n}_R[n]\\notag\\\\\n && + h_{AA}S_A[n] + n_A[n].\n \\eea\n The received signal-to-interference-plus-noise ratio (SINR) at the source ${\\tt A}$, denoted as $\\gamma_A$, is expressed as\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:hi1}\n \\gamma_A &=& \\frac{p_B|{\\bf h}_{RA}^\\dag{\\bf W} {\\bf h}_{BR}|^2}{ \\|{\\bf h}_{RA}^\\dag{\\bf W}\\|^2\n + p_A|h_{AA}|^2+ 1}.\n \\eea\n Similarly, the received SINR $\\gamma_B$ at the source ${\\tt B}$ can be\n written as\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:hi2}\n \\gamma_B &=& \\frac{p_A|{\\bf h}_{RB}^\\dag{\\bf W} {\\bf h}_{AR}|^2}{ \\|{\\bf h}_{RB}^\\dag{\\bf W}\\|^2\n + p_B|h_{BB}|^2+ 1}.\n \\eea\n The achievable rates are then given by $R_A = \\log_2(1+\\gamma_A)$ and $R_B =\n \\log_2(1+\\gamma_B)$, respectively.\n\n\\subsection{Problems Statement}\nThe conventional physical layer analog network coding scheme\nrequires two phases for ${\\tt A}$ and ${\\tt B}$ to exchange information\n\\cite{Zhang-2Phase}. Thanks to the FD operation, the proposed\nscheme reduces the whole communication to only one phase thus\nsubstantially increases the spectrum efficiency. However, the FD\noperation also brings the SI to each node so they may not always\nuse their maximum power because higher transmit power also increases\nthe level of the residual SI, therefore each node needs to carefully\nchoose its transmit power.\n\n\nWe are interested in two e2e objectives subject to each source's\npower constraints by optimizing the relay beamforming and power\nallocation at each source. The first one is to find the achievable\nrate region $(R_A, R_B)$. This can be achieved by maximizing source\n${\\tt A}$'s rate while varying the constraint on the minimum rate of\nsource ${\\tt B}$'s rate (or vice versa), i.e.,\n solving the rate maximization problem $\\mathbb{P}_1$ below:\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:prob1}\n \\mathbb{P}_1: \\max_{{\\bf W}, p_A, p_B, p_R}{R_A} ~~ \\mbox{s.t.}~~ R_B\\ge r_B, p_X\\le P_X,\n X\\in\\{A,B,R\\},\\notag\n \\eea\n where $r_B$ is the constraint on source ${\\tt B}$'s rate. By enumerating\n $r_B$, we can find the boundary of the achievable rate region.\n\nThe second problem is to maximize the sum rate of the two sources.\nMathematically, this problem is formulated as $\\mathbb{P}_2$ below:\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:prob2}\n \\mathbb{P}_2: \\max_{{\\bf W}, p_A, p_B, p_R}{R_A+R_B} ~~ \\mbox{s.t.}~~ p_X\\le P_X,\n X\\in\\{A,B,R\\}.\\notag\n \\eea\n The next two sections will be devoted to solving $\\mathbb{P}_1$ and\n $\\mathbb{P}_2$, respectively.\n\n\n\n\\section{Finding the achievable rate region}\nIn this section, we aim to optimize the relay beamforming matrix\n${\\bf W}$ and the sources' transmit power $(p_A, p_B)$ to find the\nachievable rate region. This can be achieved by solving\n$\\mathbb{P}_1$. Using the monotonicity between the SINR and the\nrate, it can be expanded as {\\small\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:fd2}\n && \\max_{{\\bf W}, p_A, p_B} ~~ \\frac{p_B|{\\bf h}_{RA}^\\dag{\\bf W} {\\bf h}_{BR}|^2}{ \\|{\\bf h}_{RA}^\\dag{\\bf W}\\|^2\n + p_A|h_{AA}|^2+ 1}\\\\\n && \\mbox{s.t.} ~~ \\frac{p_A|{\\bf h}_{RB}^\\dag{\\bf W} {\\bf h}_{AR}|^2}{ \\|{\\bf h}_{RB}^\\dag{\\bf W}\\|^2\n + p_B|h_{BB}|^2+ 1} \\ge \\Gamma_B,\\label{cons1}\\\\\n && p_A\\|{\\bf W} {\\bf h}_{AR}\\|^2 + p_B\\|{\\bf W} {\\bf h}_{BR}\\|^2 +\\textsf{trace}({\\bf W}\\qW^\\dag) \\le P_R,\\label{cons2}\\\\\n && {\\bf W}{\\bf H}_{RR}{\\bf W}={\\bf 0}, \\notag\\\\\n && 0\\le p_A\\le P_A, 0\\le p_B\\le P_B,\\notag\n \\eea}\n {\\bl where $\\Gamma_B\\triangleq 2^{r_B}-1$ is the equivalent SINR constraint for the source ${\\tt B}$.}\n Observe that all terms in \\eqref{eqn:fd2} are quadratic in ${\\bf W}$\n except the ZF constraint ${\\bf W}{\\bf H}_{RR}{\\bf W}={\\bf 0}$, which is\n difficult to handle. Considering the fact that each source only transmits a single data\n stream and the network coding principle encourages mixing rather\n than separating the data streams from different sources, we\n decompose ${\\bf W}$ as ${\\bf W} = {\\bf w}_t{\\bf w}_r^\\dag$, where\n${\\bf w}_r$ is the receive beamforming vector and ${\\bf w}_t$ is the\ntransmit beamforming vector at ${\\tt R}$. Then the ZF condition is\nsimplified to $({\\bf w}_r^\\dag{\\bf H}_{RR}{\\bf w}_t) {\\bf W} ={\\bf 0}$ or\nequivalently ${\\bf w}_r^\\dag{\\bf H}_{RR}{\\bf w}_t=0$ because in general ${\\bf W}\\ne\n{\\bf 0}$. Without loss of optimality, we further assume\n$\\|{\\bf w}_r\\|=1$. As a result, the problem \\eqref{eqn:fd2} can be\nsimplified to\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:fd711}\n \\max_{{\\bf w}_r, {\\bf w}_t, p_A, p_B} && \\frac{p_B|{\\bf h}_{RA}^\\dag {\\bf w}_t|^2|{\\bf w}_r^\\dag {\\bf h}_{BR}|^2}{ |{\\bf h}_{RA}^\\dag {\\bf w}_t|^2 + p_A|h_{AA}|^2 + 1} \\\\\n \\mbox{s.t.} && \\frac{p_A|{\\bf h}_{RB}^\\dag {\\bf w}_t|^2 |{\\bf w}_r^\\dag {\\bf h}_{AR}|^2}{ |{\\bf h}_{RB}^\\dag {\\bf w}_t|^2 + p_B|h_{BB}|^2+\n 1} \\ge \\Gamma_B,\\notag\\\\\n && p_A\\|{\\bf w}_t\\|^2 |{\\bf w}_r^\\dag{\\bf h}_{AR}|^2 + p_B\\|{\\bf w}_t\\|^2 |{\\bf w}_r^\\dag {\\bf h}_{BR}|^2 \\notag\\\\ &&+ \\|{\\bf w}_t\\|^2 \\le P_R,\\notag\\\\\n && {\\bf w}_r^\\dag{\\bf H}_{RR}{\\bf w}_t=0, \\notag\\\\\n && \\|{\\bf w}_r\\|=1,\\notag\\\\\n && 0\\le p_A\\le P_A, 0\\le p_B\\le P_B.\\notag\n \\eea\n Note that in order to guarantee the feasibility of the ZF\n constraint ${\\bf w}_r^\\dag{\\bf H}_{RR}{\\bf w}_t=0$, ${\\tt R}$ only needs to have\n two or more either transmit or receive antennas but not necessarily both.\n\n\n The problem \\eqref{eqn:fd711} is still quite complicated as\n variables ${\\bf w}_t$, ${\\bf w}_r$ and $(p_A, p_B)$ are coupled. Our idea\n to tackle this difficulty is to use an alternating optimization\n approach, i.e., at each iteration, we optimize one variable while\n keeping the other fixed, together with 1-D search to find ${\\bf w}_r$. Details are given\n below.\n\\subsection{Parameterization of the receive beamforming vector ${\\bf w}_r$}\n Observe that ${\\bf w}_r$ is mainly involved in $|{\\bf w}_r^\\dag\n{\\bf h}_{BR}|^2$ and\n $|{\\bf w}_r^\\dag {\\bf h}_{AR}|^2$, so it has to balance the signals received\n from the\n two sources. According to the result in \\cite{Jorswieck2008f},\n ${\\bf w}_r$ can be parameterized by $0\\le \\alpha\\le 1$ as below:\n\\begin{equation}} \\newcommand{\\ee}{\\end{equation}\\label{eqn:wr} {\\bf w}_r =\n\\alpha\\frac{\\Pi_{{\\bf h}_{BR}}{\\bf h}_{AR}}{\\|\\Pi_{{\\bf h}_{BR}}{\\bf h}_{AR}\\|} +\n \\sqrt{1-\\alpha}\\frac{\\Pi^\\bot_{{\\bf h}_{BR}}{\\bf h}_{AR}}{\\|\\Pi^\\bot_{{\\bf h}_{BR}}{\\bf h}_{AR}\\|}.\n \\ee\n {\\bl We have to remark that \\eqref{eqn:wr} is not the complete characterization of ${\\bf w}_r$ because it is also involved in the ZF constraint\n ${\\bf w}_r^\\dag{\\bf H}_{RR}{\\bf w}_t=0$, but this parameterization makes the\n problem more tractable.}\n\n\n Given $\\alpha$, we can optimize ${\\bf w}_t$ and $(p_A,p_B)$ as will be introduced below. Then we perform 1-D search to find the\n optimal $\\alpha^*$. We will focus on how to separately optimize\n ${\\bf w}_t$ and $(p_A, p_B)$ in the following\n two subsections. For the optimization of each variable, we will derive analytical\n solutions without using iterative methods or numerical algorithms.\n\n\n\n\n\\subsection{Optimization of the transmit beamforming vector ${\\bf w}_t$}\n We first look into the optimization of ${\\bf w}_t$ assuming ${\\bf w}_r$ and $(p_A,\n p_B)$ are fixed. Based on the problem \\eqref{eqn:fd711}, we get the\n following formulation:\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:fd71123}\n \\max_{{\\bf w}_t} && \\frac{p_B|{\\bf h}_{RA}^\\dag {\\bf w}_t|^2|{\\bf w}_r^\\dag {\\bf h}_{BR}|^2}{ |{\\bf h}_{RA}^\\dag {\\bf w}_t|^2 + p_A|h_{AA}|^2 + 1} \\\\\n \\mbox{s.t.} && \\frac{p_A|{\\bf h}_{RB}^\\dag {\\bf w}_t|^2 |{\\bf w}_r^\\dag {\\bf h}_{AR}|^2}{ |{\\bf h}_{RB}^\\dag {\\bf w}_t|^2 + p_B|h_{BB}|^2+\n 1} \\ge\\Gamma_B,\\notag\\\\\n && p_A\\|{\\bf w}_t\\|^2 |{\\bf w}_r^\\dag{\\bf h}_{AR}|^2 + p_B\\|{\\bf w}_t\\|^2 |{\\bf w}_r^\\dag {\\bf h}_{BR}|^2 \\notag\\\\ &&\n + \\|{\\bf w}_t\\|^2 \\le P_R,\\notag\\\\\n \n && {\\bf w}_r^\\dag{\\bf H}_{RR}{\\bf w}_t=0. \\notag\n \\eea\n By separating the variable ${\\bf w}_t$ and using monotonicity, \\eqref{eqn:fd71123} is\n simplified to\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:fd712}\n \\max_{{\\bf w}_t} && |{\\bf h}_{RA}^\\dag {\\bf w}_t|^2 \\\\\n \\mbox{s.t.}\n && |{\\bf h}_{RB}^\\dag {\\bf w}_t|^2 \\ge\\frac{\\Gamma_B{p_B|h_{BB}|^2+1}}{p_A|{\\bf w}_r^\\dag {\\bf h}_{AR}|^2\n -\\Gamma_B}\\triangleq\\bar\\Gamma_B,\n \\notag\\\\\n && \\|{\\bf w}_t\\|^2 \\le \\frac{P_R}{(p_A|{\\bf w}_r^\\dag{\\bf h}_{AR}|^2 + p_B |{\\bf w}_r^\\dag {\\bf h}_{BR}|^2 +1)}\\triangleq\\bar P,\\notag\\\\\n \n && {\\bf w}_r^\\dag{\\bf H}_{RR}{\\bf w}_t=0. \\notag\n \\eea\n The problem \\eqref{eqn:fd712} is not convex but it is a\n quadratic problem in ${\\bf w}_t$.\n By defining ${\\bf W}_t = {\\bf w}_t\n {\\bf w}_t^\\dag$ and using semidefinite programming relaxation, \\eqref{eqn:fd712} will become a convex problem in ${\\bf W}_t $, from which\n the optimal ${\\bf w}_t$ can be found from matrix decomposition. Interested readers are referred to \\cite{SDPR} for details. However, the special structure of the problem\n \\eqref{eqn:fd712} allows us to derive the analytical solution in the following steps.\n\\begin{enumerate}\n \\item If $\\bar\\Gamma_B<0$ or $p_A|{\\bf w}_r^\\dag {\\bf h}_{AR}|^2 <\\Gamma_B$, the\n problem is infeasible; otherwise continue.\n \\item Define the null space of the vector ${\\bf w}_r^\\dag{\\bf H}_{RR}$\n as $\\mathbf{N}_t\\in \\mathbb{C}^{M_t \\times (M_t-1)}$, i.e., ${\\bf w}_r^\\dag{\\bf H}_{RR}\\mathbf{N}_t={\\bf 0}$. Introduce a new variable ${\\bf v}\\in \\mathbb{C}^{(M_t-1)\\times 1}$ and express ${\\bf w}_t = \\mathbf{N}_t{\\bf v}$, then\n we can remove the ZF constraint in \\eqref{eqn:fd712}, and obtain\n the following equivalent problem:\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:fd712e}\n \\max_{{\\bf v}} && |{\\bf h}_{RA}^\\dag \\mathbf{N}_t{\\bf v}|^2 \\\\\n \\mbox{s.t.}\n && |{\\bf h}_{RB}^\\dag \\mathbf{N}_t{\\bf v}|^2 \\ge \\bar\\Gamma_B\n \\label{eqn:constr1} \\\\\n && \\|{\\bf v}\\|^2 \\le \\bar P, \\notag\n \\eea\n where we have used the property that $\\mathbf{N}_t^\\dag\n \\mathbf{N}_t={\\bf I}$. If $\\bar P\\|{\\bf h}_{RB}^\\dag \\mathbf{N}_t\\|^2< \\bar\\Gamma_B$, the problem\n is infeasible; otherwise continue.\n\n \\item In this step we aim to find the closed-form solution for \\eqref{eqn:fd712e}. We first solve it without the constraint\n \\eqref{eqn:constr1}.\n It can be seen that the last power constraint should always be\n satisfied with equality, and the optimal solution is given by\n ${\\bf v}^* = \\sqrt{\\bar P}\\frac{\\mathbf{N}_t^\\dag {\\bf h}_{RA}}{\\| \\mathbf{N}_t^\\dag\n {\\bf h}_{RA}\\|}$. If it also satisfies the constraint \\eqref{eqn:constr1},\n then it is the optimal solution; otherwise continue.\n \\item It this step, we know that both constraints in \\eqref{eqn:fd712e} should be active, so we reach the\n problem below:\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:fd713e}\n \\max_{{\\bf v}} && |{\\bf h}_{RA}^\\dag \\mathbf{N}_t{\\bf v}|^2 \\\\\n \\mbox{s.t.}\n && |{\\bf h}_{RB}^\\dag \\mathbf{N}_t{\\bf v}|^2 =\\bar\\Gamma_B\n \\label{eqn:constr2}\\notag\\\\\n && \\|{\\bf v}\\|^2 =\\bar P. \\notag\n \\eea\n\\end{enumerate}\n If we define ${\\bf d}_2 \\triangleq\\frac{\\mathbf{N}_t^\\dag{\\bf h}_{RA}}{\\mathbf{N}_t^\\dag{\\bf h}_{RA}}, {\\bf d}_1\\triangleq\n \\frac{\\mathbf{N}_t^\\dag{\\bf h}_{RB}}{\\mathbf{N}_t^\\dag{\\bf h}_{RB}}$, $\\phi\\in (-\\pi,\\pi]$ be the argument of\n ${\\bf d}_2^\\dag{\\bf d}_1$, $r\\triangleq|{\\bf d}_2^\\dag{\\bf d}_1|$, $q\\triangleq \\frac{\\bar\\Gamma_B}{\\bar P\\|\\mathbf{N}_t^\\dag{\\bf h}_{RB}\\|^2}$ and ${\\bf z}\\triangleq\\frac{{\\bf v}}{\\|{\\bf v}\\|}$, then we\n have the following formulation:\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:lem5}\n \\max_{{\\bf z}}&& {\\bf z}^\\dag {\\bf d}_2{\\bf d}_2^\\dag{\\bf z}, \\\\\n \\mathrm{s.t.}&& {\\bf z}^\\dag {\\bf d}_1{\\bf d}_1^\\dag{\\bf z} = q, \\quad \\|{\\bf z}\\|=1.\\notag\n \\eea\n\n The optimal solution ${\\bf z}^*$ follows from Lemma 2 in\n\\cite{Li-11} and is given below \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}{\\bf z}^*=\n \\left( r\\frac{1-q}{1-r^2}-\\sqrt{q}\\right) e^{j(\\pi-\\phi)}\n {\\bf d}_1+ \\sqrt{\\frac{1-q}{1-r^2}}{\\bf d}_2.\n\\eea\n Once we obtain ${\\bf z}^*$, the optimal transmit beamforming vector is\n given by\n \\begin{equation}} \\newcommand{\\ee}{\\end{equation} {\\bf w}_t^* = \\mathbf{N}_t {\\bf v}^* = \\sqrt{\\bar P}\\mathbf{N}_t{\\bf z}^*. \\ee\n\nAfter obtaining the optimal ${\\bf w}_t^*$ using the above procedures, we\ncan move on to find the optimal power allocation at the sources.\n\n\n\\subsection{Optimization of the source power $(p_A, p_B)$}\n Because of the FD operation at the sources and the fact that each source has a single transmit and receive antenna, they cannot suppress the residual SI in the spatial domain\n therefore cannot always use the full power. In contrast, ${\\tt R}$ has at least two transmit or receive antennas, so it can complete eliminate\n the SI and transmit using full power $P_R$. Here we aim to find the optimal power allocation $(p_A, p_B)$ at ${\\tt A}$ and ${\\tt B}$ assuming both ${\\bf w}_t$ and ${\\bf w}_r$ are\n fixed.\n\n For convenience, define $C_{At}\\triangleq |{\\bf h}_{RA}^\\dag {\\bf w}_t|^2,\nC_{rB}\\triangleq |{\\bf w}_r^\\dag {\\bf h}_{BR}|^2, C_{Bt}\\triangleq\n|{\\bf h}_{RB}^\\dag {\\bf w}_t|^2, C_{rA}\\triangleq |{\\bf w}_r^\\dag {\\bf h}_{AR}|^2\n$, then \\eqref{eqn:fd711} becomes{\\small\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:fd77}\n \\max_{p_A, p_B} && \\frac{p_B C_{At}C_{rB}}{ C_{At} + p_A|h_{AA}|^2+ 1} \\\\\n \\mbox{s.t.} && \\frac{p_A C_{Bt} C_{rA}}{ C_{Bt} + p_B|h_{BB}|^2+\n 1} \\ge \\Gamma_B\\label{eqn:linear:SINR}\\\\\n && p_A\\|{\\bf w}_t\\|^2 C_{rA} + p_B\\|{\\bf w}_t\\|^2 C_{rB} + \\|{\\bf w}_t\\|^2\\le P_R,\\label{eqn:linear:power}\\\\\n && 0\\le p_A\\le P_A, 0\\le p_B\\le P_B.\\notag\n \\eea}\nThe problem \\eqref{eqn:fd77} is a linear-fractional programming\nproblem, and can be converted to a linear programming problem\n\\cite[p. 151]{Boyd}. Again, thanks to its special structure, we can\nderive its analytical solutions below step by step.\n\\begin{enumerate}\n \\item First we check whether the constraint \\eqref{eqn:linear:SINR} is feasible. If $P_A C_{Bt} C_{rA} \\le \\Gamma_B$, then the problem is\n infeasible; otherwise, continue.\n\n \\item Next we solve \\eqref{eqn:fd77} by ignoring the constraint\n \\eqref{eqn:linear:power}. It is easy to check that at the\n optimum, at least one source should achieve its maximum power.\n The power allocation depends on two cases:\n \\begin{enumerate}\n \\item If $ P_A \\ge \\frac{\\Gamma_B( C_{Bt} \\|{\\bf w}_r\\|^2 + P_B|h_{BB}|^2+\n 1)}{C_{Bt} C_{rA} } $, then $p_B=P_B$, $p_A = \\frac{\\Gamma_B( C_{Bt} \\|{\\bf w}_r\\|^2 + P_B|h_{BB}|^2+\n 1)}{C_{Bt} C_{rA} }$; otherwise,\n \\item $p_A=P_A, p_B = \\min\\left(P_B, \\frac{\\frac{p_A C_{Bt} C_{rA}}{\\Gamma_B} - 1 - C_{Bt}\n }{|h_{BB}|^2}\\right)$.\n \\end{enumerate}\n\n \\item We check whether the above obtained solution satisfies the constraint \\eqref{eqn:linear:SINR}. If it does, then it is the optimal solution. Otherwise\n the constraint \\eqref{eqn:linear:SINR} should be met with\n equality.\n\n \\item The optimal power allocation is determined by the equation\n set below,\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n \\left\\{ \\begin{array}{l}\n p_A C_{Bt} C_{rA} =\\Gamma_B( C_{Bt} + p_B|h_{BB}|^2+\n 1), \\\\\n p_A\\|{\\bf w}_t\\|^2 C_{rA} + p_B\\|{\\bf w}_t\\|^2 C_{rB} + \\|{\\bf w}_t\\|^2\n =P_R\n \\end{array}\\right.\n \\eea\n and the solution is given by\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n \\left\\{ \\begin{array}{l}\n p_A= \\frac{\\Gamma_B|h_{BB}|^2}{C_{Bt} + C_{rA}} p_B + \\frac{\\Gamma_B(C_{Bt}+1)}{C_{Bt} +\n C_{rA}},\\\\\n p_B = \\frac{\\frac{P_R}{\\|{\\bf w}_t\\|^2} -1 - \\frac{\\Gamma_B C_{rA}(C_{Bt}+1)\n }{C_{Bt}C_{rA}}}{\\frac{\\Gamma_B C_{rA}|h_{BB}|^2\n }{C_{Bt}C_{rA}} + C_{rB} }.\n \\end{array}\\right.\n \\eea\n\\end{enumerate}\n\n\\subsection{The overall algorithm}\n Given an $\\alpha$ or ${\\bf w}_r$, we can iteratively optimize ${\\bf w}_t$ and $(p_A,\n p_B)$ as above until convergence. The value of the objective function monotonically increases as the iteration goes thus converges to a local\n optimum. We can then conduct 1-D search over $0\\le \\alpha\\le 1$ to find the\n optimal $\\alpha^*$ or ${\\bf w}_r$. By enumerating source ${\\tt B}$'s requirement $r_B$, we can numerically\nfind the boundary of the achievable rate region.\n\n {\\bl About the complexity, we remark that at each iteration, the solutions of ${\\bf w}_t$ and $(p_A,\n p_B)$ are given in simple closed-forms, so the associated complexity is low.}\n\n\n\\section{Maximizing the Sum Rate}\n In this section, we aim to maximize the sum rate of the proposed FD TWRC, i.e., to solve $\\mathbb{P}_2$, which is rewritten\n below,\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{prob:sum:rate:max}\n \\max_{{\\bf w}_t, {\\bf w}_r p_A, p_B} && \\log_2\\left(1+ \\frac{p_B C_{rB} |{\\bf h}_{RA}^\\dag {\\bf w}_t|^2}{ |{\\bf h}_{RA}^\\dag {\\bf w}_t|^2 + p_A|h_{AA}|^2 +\n 1}\\right)\\\\ &&\n +\\log_2\\left(1+ \\frac{p_A C_{rA}|{\\bf h}_{RB}^\\dag {\\bf w}_t|^2 }{ |{\\bf h}_{RB}^\\dag {\\bf w}_t|^2 + p_B|h_{BB}|^2+\n 1} \\right)\\notag \\\\\n \\mbox{s.t.} && \\|{\\bf w}_t\\|^2 \\le \\frac{P_R}{ p_A C_{rA} + p_BC_{rB} + 1},\\notag\\\\\n && {\\bf w}_r^\\dag{\\bf H}_{RR}{\\bf w}_t=0. \\notag\n \\eea\nWe will use the same characterization of \\eqref{eqn:wr} to find the optimal ${\\bf w}_r$ via 1-D\n search. We then concentrate on alternatingly optimizing the transmit beamforming vector ${\\bf w}_t$ and the power\n allocation $(p_A, p_B)$.\n \\subsection{Optimization of the transmit beamforming vector ${\\bf w}_t$}\n We first study how to optimize ${\\bf w}_t$ given ${\\bf w}_r$ and $(p_A, p_B)$.\n For convenience, we define a semidefinite matrix ${\\bf W}_t={\\bf w}_t{\\bf w}_t^\\dag$. Then the problem \\eqref{prob:sum:rate:max} becomes\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{prob:sum:rate:max2}\n \\max_{ {\\bf W}_t\\succeq {\\bf 0} } && F({\\bf W}_t) \\\\\n \\mbox{s.t.} &&\\textsf{trace}({\\bf W}) \\le \\frac{P_R}{ p_A C_{rA} + p_BC_{rB} + 1},\\notag\\\\\n && \\textsf{trace}({\\bf W}_t {\\bf H}_{RR}^\\dag{\\bf w}_r{\\bf w}_r^\\dag{\\bf H}_{RR})=0,\n \\notag\\\\\n && \\mbox{rank}({\\bf W}_t)=1, \\notag\n \\eea\n where $F({\\bf W}_t) \\triangleq \\log_2\\left(1+ \\frac{p_B C_{rB} \\textsf{trace}({\\bf W}_t{\\bf h}_{RA}{\\bf h}_{RA}^\\dag)}{ \\textsf{trace}({\\bf W}_t{\\bf h}_{RA}{\\bf h}_{RA}^\\dag) + p_A|h_{AA}|^2 + 1}\\right)\n +\\log_2\\left(1+ \\frac{p_A C_{rA} \\textsf{trace}({\\bf W}_t{\\bf h}_{RB}{\\bf h}_{RB}^\\dag) }{ \\textsf{trace}({\\bf W}_t{\\bf h}_{RB}{\\bf h}_{RB}^\\dag) + p_B|h_{BB}|^2+\n 1} \\right).$ Clearly $F({\\bf W}_t)$ is not a concave function thus (\\ref{prob:sum:rate:max2}) is a cumbersome optimization problem.\n To tackle it, we propose to\nuse the DC programming \\cite{An} to find a local optimum point. To\nthis end, we express $F({\\bf W}_t)$, as a difference of two concave\nfunctions $f({\\bf W}_t)$ and $g({\\bf W}_t)$, i.e., { \\begin{align}\n & F({\\bf W}_t)\\notag\\\\\n &= \\log_2\\left( (p_B C_{rB}+1) \\textsf{trace}({\\bf W}_t{\\bf h}_{RA}{\\bf h}_{RA}^\\dag) +\np_A|h_{AA}|^2 + 1)\n \\right) \\notag\\\\\n & - \\log_2\\left(\\textsf{trace}({\\bf W}_t{\\bf h}_{RA}{\\bf h}_{RA}^\\dag) + p_A|h_{AA}|^2 + 1\\right)\\notag\\\\\n & + \\log_2\\left( (p_A C_{rA}+1) \\textsf{trace}({\\bf W}_t{\\bf h}_{RB}{\\bf h}_{RB}^\\dag) + p_B|h_{BB}|^2 + 1)\n \\right)\\notag\\\\\n & - \\log_2\\left(\\textsf{trace}({\\bf W}_t{\\bf h}_{RB}{\\bf h}_{RB}^\\dag) + p_B|h_{BB}|^2 +\n 1\\right)\\notag\\\\\n&\\triangleq f({\\bf W}_t) - g({\\bf W}_t) \\end{align}}\n where {\\small \\begin{align} f({\\bf W}_t)\n&\\triangleq \\log_2\\left( (p_B C_{rB}+1)\n\\textsf{trace}({\\bf W}_t{\\bf h}_{RA}{\\bf h}_{RA}^\\dag) + p_A|h_{AA}|^2 + 1)\n \\right) \\notag\\\\\n & +\\log_2\\left( (p_A C_{rA}+1) \\textsf{trace}({\\bf W}_t{\\bf h}_{RB}{\\bf h}_{RB}^\\dag) + p_B|h_{BB}|^2 + 1)\n \\right),\\notag\\\\\ng({\\bf W}_t) &\\triangleq \\log_2\\left(\\textsf{trace}({\\bf W}_t{\\bf h}_{RA}{\\bf h}_{RA}^\\dag) + p_A|h_{AA}|^2 + 1\\right)\\notag\\\\\n& +\\log_2\\left(\\textsf{trace}({\\bf W}_t{\\bf h}_{RB}{\\bf h}_{RB}^\\dag) + p_B|h_{BB}|^2 +\n 1\\right).\\notag \\end{align}}\n $f({\\bf W}_t)$ is a concave function while $g({\\bf W}_t)$ is a convex\n function. The main idea is to approximate $g({\\bf W}_t)$ by a linear\n function. The linearization (first-order approximation) of $g({\\bf W}_t) $ around the point ${\\bf W}_{t,k}$\n is given by{\\small\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray} &&g_L({\\bf W}_t;{\\bf W}_{t,k}) =\n \\frac{1}{\\ln(2)}\\frac{\\textsf{trace}\\left(({\\bf W}_{t}-{\\bf W}_{t,k}){\\bf h}_{RA}{\\bf h}_{RA}^\\dag\\right)}{\\textsf{trace}({\\bf W}_{t,k}{\\bf h}_{RA}{\\bf h}_{RA}^\\dag) + p_A|h_{AA}|^2 +\n 1} \\notag\\\\\n &&+ \\frac{1}{\\ln(2)}\\frac{\\textsf{trace}\\left(({\\bf W}_{t}-{\\bf W}_{t,k}){\\bf h}_{RB}{\\bf h}_{RB}^\\dag\\right)}{\\textsf{trace}({\\bf W}_{t,k}{\\bf h}_{RB}{\\bf h}_{RB}^\\dag) + p_B|h_{BB}|^2 +\n 1} \\notag\\\\\n &&+\\log_2\\left(\\textsf{trace}({\\bf W}_{t,k}{\\bf h}_{RA}{\\bf h}_{RA}^\\dag) + p_A|h_{AA}|^2 +\n 1\\right)\\notag\\\\\n &&+\\log_2\\left(\\textsf{trace}({\\bf W}_{t,k}{\\bf h}_{RA}{\\bf h}_{RA}^\\dag) + p_A|h_{AA}|^2 +\n 1\\right).\n\\eea} Then the DC programming is applied to sequentially solve the\nfollowing convex problem, \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n{\\bf W}_{t,k+1} &= &\\arg \\max_{{\\bf W}_t}\\ f({\\bf W}_t) - g_L({\\bf W}_t; {\\bf W}_{t,k}) \\label{DCP}\\\\\n&&\\mbox{s.t.}\\quad \\textsf{trace}({\\bf W}_t) = \\frac{P_R}{ p_A C_{rA} + p_BC_{rB} + 1},\\notag\\\\\n && \\textsf{trace}({\\bf W}_t {\\bf H}_{RR}^\\dag{\\bf w}_r{\\bf w}_r^\\dag{\\bf H}_{RR})=0\n \\notag.\n\\eea To summarize, the problem \\eqref{prob:sum:rate:max2} can be\nsolved by i) choosing an initial point ${\\bf W}_t$; and ii) for $k=0, 1,\n\\cdots$, solving (\\ref{DCP}) until the termination condition is met.\nNotice that in \\eqref{DCP} we have ignored the rank-1 constraint on\n${\\bf W}_t$. This constraint is guaranteed to be satisfied by the\nresults in \\cite[Theorem 2]{Huang-11} when $M_t>2$, therefore the\ndecomposition of ${\\bf W}_t$ leads to the optimal solution ${\\bf w}_t^*$ for\n\\eqref{prob:sum:rate:max}. When $M_t=2$, the ZF constraint in the\nproblem \\eqref{prob:sum:rate:max} can determine the direction of\n${\\bf w}_t$, i.e., ${\\bf w}_t=\\sqrt{p_t}\\mathbb{N}_t$ where $p_t$ is the\ntransmit power and $\\mathbb{N}_t\\in \\mathbb{C}^{2\\times 1}$\nrepresents the null space of ${\\bf w}_r^\\dag{\\bf H}_{RR}$. Therefore the\noptimization of ${\\bf w}_t$ reduces to optimizing a scalar variable\n$p_t$, which can be found by checking the stationary points of the\nobjective function in \\eqref{prob:sum:rate:max} and the boundary\npoint without using the DC programming. The same applies to the\nspecial case of $M_t=1$.\n\n\n\\subsection{Optimization of source power $(p_A, p_B)$}\nWith ${\\bf w}_t$ and ${\\bf w}_r$ fixed, the sum rate maximization problem\n\\eqref{prob:sum:rate:max} about power allocation can be written as\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:fd777}\n \\max_{p_A, p_B} && \\log_2 \\left(1+ \\frac{p_B C_{At}C_{rB}}{ C_{At} + p_A|h_{AA}|^2+ 1}\\right)\\notag\\\\\n &&+\\log_2 \\left( 1+\\frac{p_A C_{Bt} C_{rA}}{ C_{Bt} + p_B|h_{BB}|^2+\n 1} \\right) \\\\\n \\mbox{s.t.} && p_A C_{rA} + p_B C_{rB} + 1 \\le \\frac{P_R}{\\|{\\bf w}_t\\|^2},\\label{eqn:constr3}\\\\\n && 0\\le p_A\\le P_A, 0\\le p_B\\le P_B\\notag.\n \\eea\n Note that when the first relay power constraint \\eqref{eqn:constr3} is not tight, the problem is\n the same as the conventional power allocation among two interference links to maximize the sum\n rate, and the optimal power solution is known to be binary \\cite{binary-power},\n i.e., the optimal power allocation $ (p_A^*, p_B^*)$ should satisfy\n \\begin{equation}} \\newcommand{\\ee}{\\end{equation}\n (p_A^*, p_B^*)\\in\\{(0, P_B), (P_A,0), (P_A, P_B)\\}.\n \\ee\n\n Next we only focus on the case in which the constraint \\eqref{eqn:constr3} is\n active, i.e., $p_A C_{rA} + p_B C_{rB} + 1\n =\\frac{P_R}{\\|{\\bf w}_t\\|^2}$.\nWe then have\n \\begin{equation}} \\newcommand{\\ee}{\\end{equation}\n p_A = \\frac{\\frac{P_R}{\\|{\\bf w}_t\\|^2-1}-p_B C_{rB}}{C_{rA} }.\n \\ee\n Because $0\\le p_A\\le P_A$, we can obtain the feasible range $[p_B^{\\min}, p_B^{\\max}]$ for $p_B$:\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n p_B^{\\min} &=& \\max\\left(0, \\frac{ {P_R}{\\|{\\bf w}_t\\|^2}-1 - C_{rA} P_A }{ C_{rB}}\n \\right),\\notag \\\\\n p_B^{\\max} &=& \\min \\left(P_B, \\frac{{P_R}{\\|{\\bf w}_t\\|^2}-1}{ C_{rB}}\\right).\n \\eea\n\n The objective function of \\eqref{eqn:fd777} then becomes a function of $p_B$ only, i.e.,\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n y(p_B) &=& \\log_2 (C_{At} + p_A|h_{AA}|^2+ 1+ p_B C_{At}C_{rB})\\notag\\\\&& -\\log_2( C_{At} + p_A|h_{AA}|^2+ 1)\\notag\\\\\n &&+\\log_2 ( C_{Bt} + p_B|h_{BB}|^2+ 1+ p_A C_{Bt} C_{rA})\\notag\\\\&& -\\log_2( C_{Bt} + p_B|h_{BB}|^2+\n 1),\n \\eea\n and \\eqref{eqn:fd777} reduces to a one-variable optimization, i.e.,\n \\begin{equation}} \\newcommand{\\ee}{\\end{equation}\n \\max_{p_B} ~~ y(p_B)~~~~ \\mbox{s.t.}~~ p_B^{\\min}\\le p_B\\le p_B^{\\max}.\n \\ee\n\n Setting $\\frac{\\partial y(p_B)}{\\partial p_B}=0$ leads to\n {\\small \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n&& \\frac{ C_{At}C_{rB} -\\frac{ C_{rB}}{C_{rA} } |h_{AA}|^2 }{ C_{At} + \\frac{\\frac{P_R}{\\|{\\bf w}_t\\|^2-1} |h_{AA}|^2}{C_{rA}}+ 1\n+ p_B \\left( C_{At}C_{rB} -\\frac{ C_{rB}}{C_{rA} } |h_{AA}|^2\\right) } \\\\\n&& + \\frac{\\frac{ C_{rB}}{C_{rA} } |h_{AA}|^2}{ C_{At} +\n\\frac{\\frac{P_R}{\\|{\\bf w}_t\\|^2-1}}{C_{rA} } |h_{AA}|^2 +1 -\np_B\\frac{ C_{rB}}{C_{rA} } |h_{AA}|^2 }\\notag\\\\\n&& +\\frac{ |h_{BB}|^2- \\frac{ C_{rB}}{C_{rA} } C_{Bt}\n C_{rA} }{ C_{Bt} + \\frac{\\frac{P_R}{\\|{\\bf w}_t\\|^2-1}}{C_{rA} } C_{Bt} C_{rA}+ 1+ p_B\\left( |h_{BB}|^2- \\frac{ C_{rB}}{C_{rA} } C_{Bt}\n C_{rA}\\right)}\\notag\\\\\n&& - \\frac{|h_{BB}|^2}{ C_{Bt} + p_B|h_{BB}|^2+1}\n =0.\\notag\n \\eea}\nThis in turn becomes a cubic (3-rd order) equation and all roots\ncan be found analytically. Suppose that the set of all positive root\nwithin $(p_B^{\\min}, p_B^{\\max})$ is denoted as $\\Psi$ which may\ncontain 0, 1 or 3 elements. In order to find the optimal $p_B^*$,\nwe need to compare the objective values of all elements in the set\n$\\Psi\\cup\\{p_B^{\\min}, p_B^{\\max}\\}$ and choose the one that results\nin the maximum objective value.\n\n\n{\\bl We comment that the complexity of the overall algorithm to\nmaximize the sum rate is dominated by the optimization of ${\\bf w}_t$\nand more specially, solving the problem \\eqref{DCP}. Since\n\\eqref{DCP} is a semi-definite programming (SDP) problem with one\nvariable ${\\bf W}_t\\in \\mathbb{C}^{M_t\\times M_t}$ and two constraints,\nthe worst-case complexity to solve it is\n$\\mathcal{O}(M_t^{4.5}\\log(\\frac{1}{\\epsilon}))$ where $\\epsilon$ is\nthe desired solution accuracy \\cite{SDPR}.}\n\n\n\\section{Benchmark schemes}\nIn this section, we introduce three benchmark schemes that the\nproposed FD network scheme can be compared with. The first one is\nthe conventional two-phase HD TWRC using the analog network coding,\nwhich is known to outperform the three-phase and four-phase HD\nschemes to provide throughput gain\\cite{practical_PLNC}; the second\none is a two-phase one-way FD scheme and in each phase, the relay\nworks in the FD mode\\cite{FD-two-phase}; the last one ignores the\nresidual SI channel at the relay thus provides a performance upper\nbound that is useful to evaluate the proposed algorithms.\n\n\\subsection{Two-phase HD relaying using analog network coding}\n HD analog network coding is introduced in \\cite{Zhang-2Phase} which takes\n two phases to complete information exchange between ${\\tt A}$ and\n ${\\tt B}$. In the first phase, both sources transmit to ${\\bf R}$ and in\n the second phase, the relay multiplies the received signal by a\n beamforming matrix ${\\bf W}$ then broadcasts it to ${\\tt A}$ and\n ${\\tt B}$. Because there is no SI, every node can use its full power,\n so only ${\\bf W}$ needs to be optimized. The achievable rate pair\n and the relay power consumption, are given by\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n R_A &=& \\frac{1}{2}\\log_2\\left(1+ \\frac{P_B|{\\bf h}_{RA}^\\dag{\\bf W} {\\bf h}_{BR}|^2}{ \\|{\\bf h}_{RA}^\\dag{\\bf W}\\|^2\n +1}\\right),\\notag\\\\\n R_B &=& \\frac{1}{2}\\log_2\\left(1+ \\frac{P_A|{\\bf h}_{RB}^\\dag{\\bf W} {\\bf h}_{AR}|^2}{\n\\|{\\bf h}_{RB}^\\dag{\\bf W}\\|^2 + 1}\\right),\\notag\\\\\n p_R &=& \\|{\\bf W} {\\bf h}_{AR}\\|^2 P_A + \\|{\\bf W} {\\bf h}_{BR}\\|^2 P_B\n +\\textsf{trace}({\\bf W}\\qW^\\dag).\\notag\n \\eea\n where the factor of $\\frac{1}{2}$ is due to the two\n transmission phases used.\n Given the above rates and the power expression, problems\n $\\mathbb{P}_1$ and $\\mathbb{P}_2$ can be solved to find the\n achievable rate region and the maximum sum rate, respectively.\n Compared with this scheme, the proposed FD relaying can reduce the\n total communication phases to one, thus has the potential to\n improve the throughput.\n\n\\subsection{Two-phase one-way FD}\n Another scheme that we will compare with is the FD one-way relaying\n in which ${\\tt R}$ works in the FD mode while the two sources work in the\n HD mode. In this way, both sources can transmit with the maximum power. We use the same notation as the proposed scheme and the relay beamforming matrix is ${\\bf W}={\\bf w}_t{\\bf w}_r^\\dag$.\n\n The achievable rate, relay power constraint, and zero residual SI constraint for source ${\\tt A}$ (direction: ${\\tt B} \\rightarrow{\\tt A}$) are, respectively,\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\n &&R_A =\\log_2\\left(1+ \\frac{P_B|{\\bf h}_{RA}^\\dag {\\bf w}_t|^2|{\\bf w}_r^\\dag {\\bf h}_{BR}|^2}{ |{\\bf h}_{RA}^\\dag {\\bf w}_t|^2 + 1} \\right), \\\\\n &&p_B\\|{\\bf w}_t\\|^2 |{\\bf w}_r^\\dag {\\bf h}_{BR}|^2 +\n\\|{\\bf w}_t\\|^2\\|{\\bf w}_r\\|^2\\le P_R, \\\\\n &&{\\bf w}_r^\\dag{\\bf H}_{RR}{\\bf w}_t=0.\\label{eqn:ZF}\n \\eea\n\n\n In our previous work \\cite{MIMO-relay-HD}, we have derived the closed-form expressions below for $R_A$\n depending on how the ZF constraint is realized:\n \\begin{enumerate}\n\\item\n Receive ZF. In this case, we assume ${\\bf w}_t={\\bf h}_{RA}$ and choose\n ${\\bf w}_r$ to achieve \\eqref{eqn:ZF}. We showed that the achievable\ne2e received signal-to-noise ratio (SNR) can be expressed as\n \\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:SNR:R}\n \\gamma_{RZF}\n =\\frac{P_B \\|{\\bf D}{\\bf h}_{BR}\\|^2 P_R \\|{\\bf h}_{RA}\\|^2}\n {P_B \\|{\\bf D}{\\bf h}_{BR}\\|^2 + P_R\\|{\\bf h}_{RA}\\|^2 + 1},\n \\eea\n where ${\\bf D}\\triangleq \\Pi^\\bot_{{\\bf H}_{RR}{\\bf h}_{RA}}$.\n\n\n\\item Transmit ZF. In this case, we assume ${\\bf w}_r={\\bf h}_{BR}$ and choose\n ${\\bf w}_t$ to achieve \\eqref{eqn:ZF}. We then reach the following\n achievable SNR:\n\\begin{eqnarray}} \\newcommand{\\eea}{\\end{eqnarray}\\label{eqn:SNR:t}\n \\gamma_{TZF} = \\frac{ P_B \\|{\\bf h}_{BR}\\|^2 P_R\\|{\\bf B}{\\bf h}_{RA}\\|^2}{\n {P_B\\|{\\bf h}_{BR}\\|^2 + {P_R}\\|{\\bf B}{\\bf h}_{RS}\\|^2+ 1 }},\n\\eea where ${\\bf B}\\triangleq \\Pi_{{\\bf H}_{RR}^\\dag{\\bf h}_{BR}}^\\bot$.\n\\end{enumerate}\n $R_A$ is then determined by $R_A=\\log_2(1+\\max(\\gamma_{RZF},\n \\gamma_{TZF}))$. We can derive similar achieve rate $R_B$ for the source\n ${\\tt B}$.\n\n Note that $R_A$ and $R_B$ cannot be achieved simultaneously as it requires\n that each corresponding source occupies the whole transmission time.\n The boundary of the rate region can be obtained by using time-sharing parameter $t\\in[0,1]$, i.e., $(t R_A, (1-t)R_B)$.\n\n\\subsection{FD-Upper Bound}\n This scheme is the same as the proposed FD scheme except that we\n assume there is no SI at ${\\tt R}$, but we still consider the SI at the two sources, i.e., ${\\bf H}_{RR}={\\bf 0}, |h_{AA}|>0, |h_{BB}|>0$.\n In this case the ZF constraint in \\eqref{eqn:fd2} is not necessary.\n We remark that this scheme uses unrealistic assumption\n of ${\\bf H}_{RR}={\\bf 0}$, so it is not a practical scheme but provides a\nuseful upper bound to evaluate the performance of the proposed\nalgorithms.\n\nIn the simulation results, we will label the above three benchmark\nschemes as ``Two-phase HD'', ``Two-phase FD'' and ``Proposed\none-phase FD upper bound'', respectively.\n\n\n\n\\section{Numerical Results}\\label{sec:simu}\nIn this section we provide numerical results to illustrate the\nachievable rate region and the sum rate performance of the proposed\nFD two-way relaying scheme. We compare it with the above mentioned\nthree benchmark schemes. The simulation set-up follows the system\nmodel in Section II. Unless otherwise specified, we assume that\nthere are $M_T=M_R=3$ antennas at ${\\tt R}$, the average residual SI\nchannel gain is $\\sigma^2_A=\\sigma^2_B=\\sigma^2_R=-20$ dB, {\\bl and\nthe per-node transmit SNR for both sources and the relay is\n$P_A=P_B=P_R=10$ dB, which are the power constraints in problem\nformulations $\\mathbb{P}_1$ and $\\mathbb{P}_2$}. The results are\nobtained using 100 independent channel realizations.\n\n\n\\input figures.tex\n\n\n\\subsection{Achievable rate region}\nFirst we illustrate the achievable {\\bl average} rate region for the\ntwo sources ${\\tt A}$ and ${\\tt B}$ in Fig. \\ref{fig:rate:region}. {\\bl It\nis seen that the two-phase FD scheme already greatly enlarges the\nachievable data region of the conventional HD scheme. The proposed\none-phase FD scheme achieves significantly larger rate region over\nthe two-phase FD and the conventional HD schemes.} We observe that\nthere is still a noticeable gap between the proposed scheme and its\nupper bound because the proposed solution is suboptimal.\n\n\n\n\\subsection{Sum rate performance}\nWe then investigate the effect of the source transmit SNR $P_A$ and\n$P_B$ ($P_A=P_B$) on the sum rate shown in Fig.\n\\ref{fig:sumrate:R1020dB} when { $P_R=10$ dB (solid curves) and\n$P_R=20$ dB (dashed curves), respectively}. {\\bl We first consider\nthe case $P_R=10$ dB.} As expected, the sum rate improves as the\nsource transmit SNR increase and the proposed one-phase FD schemes\nclearly outperforms the two benchmark schemes. When the source SNR\nis above 15 dB, the sum rate of the proposed FD scheme saturates.\nThis is because the high transmit power results in high residual SI,\ntherefore increasing power budget does not necessarily improve the\nperformance. To illustrate the performance improvement in sum rate,\nwe show the sum rate gain over the conventional two-phase HD scheme\nin Fig. \\ref{fig:sumrate:gain:R1020dB}. It is observed that when\nthe source SNR is 10 dB, the proposed one-phase FD scheme and the\ntwo-phase FD scheme can achieve the sum rate gain of 1.56 and 1.22,\nrespectively. Even the performance upper bound cannot achieve double\nrate because of the residual SI at both sources. The rate gain in\ngeneral decreases as the source transmit SNR increases again because\nof the residual SI therefore the sources need to carefully adjust\nits transmit power.\n\n\nThe same trend is observed when $P_R=20$ dB. In this case, a sum\nrate gain of nearly 1.7 is recorded when the source transmit SNR is\n10 dB. Another interesting observation is that the performance of\nthe proposed scheme is very close to the upper bound. { This is\nbecause when the relay power is high, the e2e performance is limited\nby the link from the source to the relay, rather than the relay to\nthe other source, therefore the residual SI at the relay has little\neffect on the sum rate.}\n\nThe impact of the relay transmit SNR on the sum rate is shown in\nFig. \\ref{fig:sumrate:RPower}. It is seen that when the relay SNR is\nlow, the proposed one-phase FD scheme achieves lower sum rate than\nthe two-phase FD scheme at low transmit SNRs then outperforms the\nlatter when the transmit SNR is above 5 dB. The performance gain is\nremarkable when the relay transmit SNR is high. This is because\nunlike the two sources, the relay can null out the residual SI\nusing multiple antennas, therefore it can always use the maximum\navailable power to improve the sum rate.\n\nThe effect of the residual SI channel gain at the sources is\nexamined in Fig. \\ref{fig:sumrate:SI}. Naturally, the sum rate\ndecreases as the residual SI channel becomes stronger or the SI is\nnot adequately suppressed. When the residual SI channel gain is\nabove -5 dB, the two-phase FD scheme outperforms the one-phase FD\nscheme while both still achieve higher sum rate than the\nconventional two-phase HD scheme even when the SI channel gain is\nas high as 5 dB.\n\n\n Next we show the sum rate and rate gain results when the number of antennas ($M_R=M_T$) at the\nrelay varies from 2 to 6 in Fig. \\ref{fig:sumrate:antennas} and\n\\ref{fig:sumrate:gain:antennas}, respectively. The sum rate steady\nincreases as more antennas are placed at the relay due to the array\ngain. It is observed that the rate gain remains about 1.55 when the\nnumber of antennas is greater than 2.\n\n\n\\subsection{Asymmetric channel gain}\n The above results are mainly for a symmetric case, i.e., both sources have\n similar power constraints and channel {\\bl strengths}. Here we consider an\n asymmetric case where the average channel gain between ${\\tt R}$ and\n ${\\tt B}$ is -10 dB. We plot the rate region in Fig. \\ref{fig:sumrate:region:asymmetric}\nusing the same system parameters as those in Fig. 2 except that the\ngain of channel vectors ${\\bf h}_{BR}$ and\n ${\\bf h}_{RB}$ is 10 dB weaker. {\\bl The results show that both sources'\n rates are reduced while the source ${\\tt B}$ suffers more rate loss.} This is because one\n source's channels to and from the relay will also affect the\n performance of the other source. The sum rate comparison is given\n in Fig. \\ref{fig:sumrate:asymmetric} and it is observed the performance of the two-phase {\\bl FD} scheme is very close to that of the proposed FD scheme at all SNR region. This can be explained\n by the fact that e2e performance is restricted by the channel quality between ${\\tt R}$\n and ${\\tt B}$, so the gain due to the simultaneous transmission of two sources is limited.\n\n{\n\\subsection{Impact of the local channel state information}\n Finally we consider the case that only the receive CSI at each node is available but the transmit CSI is unknown. Because of the lack of the transmit\n CSI, the two sources use full power\n $P_A$ and $P_B$; ${\\bf w}_t$ at the relay is chosen arbitrarily to satisfy the ZF constraint and the relay power\n constraint. The sum rate performance is shown in Fig.\n \\ref{fig:sumrate:csi}. It is seen that the proposed FD scheme still achieves significant performance gain over the HD relaying at low to medium transmit SNRs although all rates\n are much lower than the case with the global CSI in Fig. 3. Another\n notable difference is that at high transmit SNRs, the performance of the proposed FD scheme degrades quickly.\n This is because the two sources need to adjust its transmit power\n rather than using full power. This highlights the importance of the\n global CSI for adapting the transmit power and adjusting the relay beamforming.\n}\n\n\n \\section{Conclusion and Future work}\n We have investigated the application of the FD operation to MIMO TWRC, which requires only one phase for the two sources to exchange\n information. We studied two problems of finding the achievable\n rate region and maximizing the sum rate by optimizing the relay\n beamforming matrix and power allocation at the sources.\n Iterative algorithms are proposed together with 1-D search to\n find the local optimum solutions. At each iteration, either\n analytical solution or convex formulation has been derived. We\n have conducted intensive simulations to illustrate the effects of\n different system parameters. The results show that the FD\n operation has great potential to achieve much higher data rates than the conventional HD TWRC.\n\n\n{\\bl Regarding the future directions, better suboptimal and the\noptimal solutions are worth studying.} There are a couple of reasons\nwhy the proposed algorithm is sub-optimal such as the additional ZF\nconstraint, {\\bl the incomplete characterization of the receive\nbeamforming vector}, and the alternating optimization algorithms.\nAnother direction is to study the use of multiple transmit\/receive\nantennas at the two sources. If a single data stream is transmitted,\nthe residual SI can be removed using the ZF criterion at the sources\nas well. This actually simplifies the optimization as the two\nsources can use the maximum power. However, multiple antennas can\nsupport multiple and variable number of data streams, and when the\nproblem is coupled with the SI suppression, it will be much more\nchallenging. Thirdly, in this paper, we focus on the benefit of the\nFD {\\bl in terms of} spectrum efficiency. In \\cite{practical_PLNC},\nit is shown that for the HD case, three-phase transmission schemes\noffers a better compromise between the sum rate and the bit error\nrate than the two-phase scheme, especially in the asymmetric case.\nIt is worth investigating whether such an trade-off also exists for\nthe FD scenario. \\vspace{-2mm}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nRecent observations of Type IA Supernovae \\cite{SN1,SN2} as well as\nconcordance with other observations (including the microwave\nbackground and galaxy power spectra) indicate that the universe is\naccelerating. Possible explanations for such an acceleration include\na variety of dark energy models in which the universe is vacuum\ndominated, such as a cosmological constant. In 1986, we explored the\npossibility of a time-dependent vacuum \\cite{fafm} characterized by an\nequation of state $w$. Quintessence models\n\\cite{ratpeeb,frieman,wett,stein,caldwell} utilize a rolling scalar\nfield to achieve time dependence of the vacuum.\n\nAs an alternative approach to explain the acceleration, we proposed\nmodifications to the Friedmann equation in lieu of having any vacuum\nenergy at all \\cite{freeselewis} (hereafter Paper I). In our\nCardassian model, \\footnote{The name Cardassian refers to a humanoid\n race in Star Trek whose goal is accelerated expansion of their evil\n empire (aka George W). This race looks foreign to us and yet is made\n entirely of matter.} the universe is flat and accelerating, and yet\nconsists only of matter and radiation.\nThe usual Friedmann equation governing the expansion of the universe\n\\begin{equation}\n\\label{eq:usual}\nH^2 = \\left({\\dot a \\over a}\\right)^2 = {8 \\pi G \\over 3} \\rho \n\\end{equation}\nis modified to become\n\\begin{equation}\n\\label{eq:general}\nH^2 = \\left({\\dot a \\over a}\\right)^2 = g(\\rho) ,\n\\end{equation}\nwhere $\\rho$ contains only matter and radiation (no vacuum), $H=\\dot\na\/a$ is the Hubble constant (as a function of time), $G=1\/m_{pl}^2$ is\nNewton's gravitational constant, and $a$ is the scale factor of the\nuniverse. We note here that the geometry is flat, as required by\nmeasurements of the cosmic background radiation \\cite{boom,dasi}, so that\nthere are no curvature terms in the equation. There is no vacuum term\nin the equation. The model does not address the cosmological constant\n($\\Lambda$) problem; we simply set $\\Lambda=0$.\n\nWe take $g(\\rho)$ to be a function of $\\rho$ that returns simply to $8\n\\pi G \\rho\/3$ at early epochs, but that takes a different form that\ndrives an accelerated expansion in the recent past of the universe at\n$z<{\\cal O}(1)$.\n\nI begin by describing the phenomenology of Cardassian models, and then\nturn to the motivation for modified Friedmann equations of this form.\nThe new term required may arise, e.g., as a consequence of our\nobservable universe living as a 3-dimensional brane in a higher\ndimensional universe. A second possible interpretation of Cardassian\nexpansion is developed \\cite{gondolo}, in which we treat the modified\nFriedman equations as due to a fluid, in which the energy density has\nnew contributions with negative pressure (possibly due to dark matter\nwith self-interactions).\n\n\\section{Power Law Cardassian Model}\n\nThe simplest version of Cardassian expansion invokes the addition of a\nnew power law term to the right hand side of the Friedmann equation:\n\\begin{equation}\n\\label{eq:new}\nH^2 ={8\\pi G\n\\over 3 } \\rho + B \\rho^n \n\\end{equation}\nwhere $n$ is a number with\n\\begin{equation}\nn<2\/3 .\n\\end{equation}\nThe new term is initially negligible, and only comes to dominate at\nredshift $z \\sim {\\cal O}(1)$. Once it dominates, it causes the\nuniverse to accelerate.\n\nWe take the usual energy conservation:\n\\begin{equation}\n\\label{eq:energy}\n\\dot \\rho + 3H (\\rho + p) = 0 ,\n\\end{equation}\nwhich gives the evolution of matter:\n\\begin{equation}\n\\rho_M = \\rho_{M,0}(a\/a_0)^{-3} .\n\\end{equation}\nHere subscript $0$ refers\nto today. \nEqs.(\\ref{eq:general}) and (\\ref{eq:energy})\ncontain the complete information of the expansion history.\n\nThe new term in Eq.(\\ref{eq:new}) (the second term on the right hand side)\nis initially negligible; hence ordinary early universe cosmology,\nincluding nucleosynthesis, results.. The new term only comes to\ndominate recently, at the redshift $z_{car} \\sim O(1)$ indicated by\nthe supernovae observations. Once the second term dominates, it\ncauses the universe to accelerate. When the new term is so large that\nthe ordinary first term can be neglected, we find\n\\begin{equation}\nR \\propto t^{2 \\over 3n}\n\\end{equation}\nso that the expansion is superluminal (accelerated) for $n<2\/3$. As\nexamples, for $n=2\/3$ we have $R \\sim t$; for $n=1\/3$ we have $R \\sim\nt^2$; and for $n=1\/6$ we have $R \\sim t^4$. \nFor $n<1\/3$ the acceleration is increasing (the cosmic jerk).\n\nThere are two free parameters in the power law model, $B$ and $n$ (or\nequivalently, $z_{car}$ and $n$). We choose one of these parameters\nto make the second term term kick in at the right time to explain the\nobservations. As yet we have no explanation of the coincidence\nproblem; i.e., we have no explanation for the timing of $z_{car}$.\nSuch an explanation may arise in the context of extra dimensions.\n\nObservations of the cosmic background radiation show that the geometry\nof the universe is flat with $\\Omega_0=1$. How can we reconcile this\nwith a theory made entirely of matter (and radiation) and yet\nobservations of the matter density that indicate $\\Omega_M =\n0.3$? In the Cardassian model we need to revisit the question of what\nvalue of energy density today, $\\rho_0$, corresponds to a flat\ngeometry. We will show that the energy density required to close the\nuniverse is much smaller than in a standard cosmology, so that matter\ncan be sufficient to provide a flat geometry.\n\nFrom evaluating Eq.(\\ref{eq:new}) today, we have\n\\begin{equation}\n\\label{hubbletoday}\nH_0^2 = A \\rho_0 + B \\rho_0^n .\n\\end{equation}\nThe energy density $\\rho_0$ that satisfies Eq.(\\ref{hubbletoday}) is,\nby definition, the critical density. The usual value of the critical\ndensity is found by solving the equation with only the first term\n(with $B=0$), so that\n\\begin{equation}\n\\rho_{c} = {3 H_0^2\n\\over 8 \\pi G} = 1.88 \\times 10^{-29} h_0^2 {\\rm gm\/cm^{-3}}\n\\end{equation}\nand $h_0$ is the Hubble constant today in units of 100 km\/s\/Mpc.\nHowever, in the presence of the second term, we can solve\nEq.(\\ref{hubbletoday}) to find that the critical density $\\rho_c$ has\nbeen modified from its usual value to a different number $\\tilde\n\\rho_c$. We use our second free parameter to fix this number to give\n\\begin{equation}\n\\rho_{M,0} = \\tilde \\rho_c = 0.3 \\rho_c \n\\end{equation}\nso that $\\Omega_M^{obs} = \\rho_{M,0}\/\\rho_c = 0.3$.\n\nHence the universe can be flat, matter dominated, and accelerating,\nwith a matter density that is 0.3 of the old critical density.\nMatter can provide the entire closure density of the Universe.\n\nAn equivalent formulation of power law Cardassian is\n\\begin{equation}\n\\label{eq:new2}\nH^2 = A \\rho [1 + ({\\rho \\over \\rho_{car}})^{n-1}] .\n\\end{equation}\nThe first term inside the bracket dominates initially but the second\nterm takes over once the energy density has dropped to the value\n$\\rho_{car}$. Here, $\\rho_{car}$ is the energy density at which the\ntwo terms are equal: the ordinary energy density term on the right\nhand side of the FRW equation is equal in magnitude to the new term.\nFor reasonable parameters, $\\rho_{car} \\sim 2.7 \\tilde \\rho_c \\sim\n10^{-29}$gm\/cm$^3$. Since the modifications are important only for $\\rho\n< \\rho_{card}$, solar system physics is completely unaffected.\n\n\\subsection{Observational tests of power law Cardassian}\n\nThe power law Cardassian model with an\nadditional term $\\rho^n$ satisfies many observational constraints: the\nuniverse is somewhat older, the first Doppler peak in the microwave\nbackground is slightly shifted, early structure formation ($z>1$) is\nunaffected, but structure will stop growing sooner. In addition the\nmodifications to the Poisson equation will affect cluster abundances and\nthe ISW affect in the CMB.\n\n\\subsection{Comparing to Quintessence}\n\nWe note that, with regard to observational tests, one can make a\ncorrespondence between the $\\rho^n$ Cardassian and Quintessence models\nfor constant $n$; we stress, however, that the two models are entirely\ndifferent. Quintessence requires a dark energy component with a specific\nequation of state ($p = w\\rho$), whereas the only ingredients in the\nCardassian model are ordinary matter ($p = 0$) and radiation ($p =\n1\/3$). However, as far as any observation that involves only $R(t)$, or\nequivalently $H(z)$, the two models predict the same effects on the\nobservation. Regarding such observations, we can make the following\nidentifications between the Cardassian and quintessence models: $n\n\\Rightarrow w+1$, $F\\Rightarrow \\Omega_m$, and $1-F \\Rightarrow\n\\Omega_Q$, where $w$ is the quintessence equation of state parameter,\n$\\Omega_m= \\rho_m\/\\rho_{c}$ is the ratio of matter density to the\n(old) critical density in the standard FRW cosmology appropriate to\nquintessence, $\\Omega_Q= \\rho_Q\/\\rho_{c}$ is the ratio of\nquintessence energy density to the (old) critical density, and \n$F=\\tilde{\\rho_c}\/\\rho_c$. In this way, the Cardassian model with\n$\\rho^n$ can make contact with quintessence with regard to observational\ntests. We note that Generalized Cardassian models can be distinguished\nfrom generic quintessence models with upcoming precision cosmological\nexperiments.\n\n\\subsection{Best Fit of Parameters to Current Data}\n\nWe can find the best fit of the Cardassian parameters $n$ and\n$z_{car}$ to current CMB and Supernova data. The current best fit is\nobtained for $w \\leq -0.78$, or, equivalently, $n \\leq 0.22$\n\\cite{wmap} and $0.33 \\leq z_{card} \\leq 0.48$. As an example, for\n$n= 0.2$ (equivalently, $w=-0.8$), we find that $z_{car} = 0.42$.\nThen the position of the first Doppler peak is shifted by a few\npercent. The age of the universe is 13 Gyr \\cite{savage}. There are\nsome indications that $w \\leq -1$ may give a good fit \\cite{mmot} and\nthis case of phantom cosmology will be discussed below.\n\n\\section{Generalized Cardassian}\n\nMore generally, we consider other forms of $g(\\rho)$ in Eq.(3), as\ndiscussed in \\cite{conf}. For example, a simple generalization of\nEq.(\\ref{eq:new}) is Modified Polytropic Cardassian:\n\\begin{equation}\n\\label{eq:MPC}\nH^2 = {8 \\pi G \\over 3} \\rho [1+ ({\\rho\/\\rho_{car}})^{q(n-1})]^{1\/q} \n\\end{equation}\nwith $q>0$ and $n<2\/3$. The right hand side returns to the ordinary\nFriedmann equation at early times, but becomes $\\rho^n$ at late times,\njust as in Eq.(\\ref{eq:new2}). Other examples of Generalized\nCardassian have been discussed in \\cite{conf,gondolo}.\n\n\\subsection{Observational Tests of Generalized Cardassian Models}\n\n\\subsubsection{Current Supernova Data}\n\nFigure 1 shows comparisons by Wang et al \\cite{wang} of MP Cardassian\nmodels with current supernova data. Clearly the models fit the\nexisting data. The figure also shows predictions out to $z=2$ for\nfuture supernova data. Future supernova data will thus allow one to\ndifferentiate between Cardassian models, quintessence models, and\ncosmological constant models.\n\n\\begin{figure}[ht]\n\\includegraphics[width=6in,height=6in]{SNnow.eps}\n\\caption{\\label{Figure 1}\n Comparison of MP Cardassian models (see Eq.(\\ref{eq:MPC}) with\n current Supernova data, as well as predictions out to $z=2$ for\n future data. For comparison, dotted lines indicate $(\\Omega_m,\n \\Omega_{\\Lambda}$ as labelled on the right hand side of the plot.\n We see that MP Cardassian models with $n=0.2$ and $q$ ranging from\n 1-3 match the current data, and can be differentiated from generical\n cosmological constant or quintessence models in future data. }\n\\end{figure}\n\n\\subsubsection{Additional Observational Tests}\n\nA number of additional observational tests can be used to test the\nmodels. We plan to use the code CMBFAST to further constrain\nCardassian parameters in light of Cosmic Background Radiation data\n(WMAP). In particular, the Integrated Sachs Wolfe effect may\ndifferentiate Cardassian from generic quintessence models. A second\napproach is number count tests such as DEEP2. The abundance of galaxy\nhaloes of fixed rotational speed depend on the comoving volume\nelement. Third, the Alcock-Paczynski test compares the angular size\nof a spherical object at redshift $z$ to its redhsift extent $\\Delta\nz$. Depending on the cosmology, the spherical object may look\nsquooshed in redshift space. A proposed trick \\cite{crottshui} is to\nuse the correlation function of Lyman-alpha clouds as spherical\nobjects.\n\n\\section{Motivation for Cardassian Cosmology}\n\nWe present two possible origins for Cardassian models.\nThe original idea arose from consideration of braneworld scenarios,\nin which our observable universe is a three dimensional membrane\nembedded in extra dimensions. Recently, as a second interpretation\nof Cardassian cosmology, we have investigated\nan alternative four-dimensional fluid description. Here,\nwe take the ordinary Einstein equations, but the energy density\nhas additional terms (possibly due to dark matter with self-interactions\ncharacterized by negative pressure).\n\n\\subsection{Braneworld Scenarios}\n\nIn braneworld scenarios, our observable universe is a three\ndimensional membrane embedded in extra dimensions. For simplicity, we\nwork with one extra dimension. In this five-dimensional\nworld, we take the metric to be\n\\begin{equation}\n\\label{eq:metric}\nds^2 = -q^2(\\tau,u)d\\tau^2 + a^2(\\tau,u)d\\vec{x}^2 + b^2(\\tau,u)du^2\n\\end{equation}\nwhere $u$ is the coordinate in the direction of the fifth dimension,\n$a$ is the usual scale factor of our observable universe and $b$ is\nthe scale factor in the fifth dimension. We found \\cite{cf}\nthat one does not generically obtain the usual Friedmann equation\non the observable brane. In fact, by a suitable choice of the\nbulk energy momentum tensor, one may obtain a Friedmann equation\non our brane of the form $H^2 \\sim \\rho^n$ for any $n$. More\ngenerally, one can obtain other modifications. \\cite{bin} showed\nthat for an antideSitter bulk one obtains a quadratic correction\n$H^2 \\sim \\rho^2$, but more generally one can obtain any power of $n$,\nas shown in \\cite{cf}. \n\nThe origin of these modifications to the Friedmann equations is as\nfollows. We consider the five dimensional Einstein equations,\n\\begin{equation}\n\\label{eq:5dEinstein}\n\\tilde{G}_{AB} = 8\\pi G_{(5)} \\tilde{T}_{AB} ,\n\\end{equation}\nwhere the 5D Newton's constant is related to the 5D Planck mass as\n$ 8\\pi G_{(5)} = M_{(5)}^{-3}$.\nThe 5D Einstein tensor on the left hand side is given by\n\\begin{eqnarray}\n{\\tilde G}_{00} &=& 3\\left\\{ \\frac{\\da}{a} \\left( \\frac{\\da}{a}+ \\frac{\\db}{b} \\right) - \\frac{n^2}{b^2} \n\\left(\\frac{\\ppa}{a} + \\frac{\\pa}{a} \\left( \\frac{\\pa}{a} - \\frac{\\pb}{b} \\right) \\right) \\right\\}, \n\\label{ein00} \\\\\n {\\tilde G}_{ij} &=& \n\\frac{a^2}{b^2} \\delta_{ij}\\left\\{\\frac{\\pa}{a}\n\\left(\\frac{\\pa}{a}+2\\frac{\\pn}{n}\\right)-\\frac{\\pb}{b}\\left(\\frac{\\pn}{n}+2\\frac{\\pa}{a}\\right)\n+2\\frac{\\ppa}{a}+\\frac{\\ppn}{n}\\right\\} \n\\nonumber \\\\\n& &+\\frac{a^2}{n^2} \\delta_{ij} \\left\\{ \\frac{\\da}{a} \\left(-\\frac{\\da}{a}+2\\frac{\\dn}{n}\\right)-2\\frac{\\dda}{a}\n+ \\frac{\\db}{b} \\left(-2\\frac{\\da}{a} + \\frac{\\dn}{n} \\right) - \\frac{\\ddb}{b} \\right\\},\n\\label{einij} \\\\\n{\\tilde G}_{05} &=& 3\\left(\\frac{\\pn}{n} \\frac{\\da}{a} + \\frac{\\pa}{a} \\frac{\\db}{b} - \\frac{\\dot{a}^{\\prime}}{a}\n \\right),\n\\label{ein05} \\\\\n{\\tilde G}_{55} &=& 3\\left\\{ \\frac{\\pa}{a} \\left(\\frac{\\pa}{a}+\\frac{\\pn}{n} \\right) - \\frac{b^2}{n^2} \n\\left(\\frac{\\da}{a} \\left(\\frac{\\da}{a}-\\frac{\\dn}{n} \\right) + \\frac{\\dda}{a}\\right) \\right\\}.\n\\label{ein55} \n\\end{eqnarray} \nIn the above expressions, a prime stands for a derivative with respect to\n $u$, and a \ndot for a derivative with respect to $\\tau$. \n\nThese 5D Einstein equations are supplemented by boundary conditions\n(known as Israel conditions) due to the existence of our 3-brane.\nThese boundary conditions are exactly analogous to those of a charged\nplate in electromagnetism. There, the change in the perpendicular\ncomponent of the electric field as one crosses a plate with charge per\narea $\\sigma$ is given by $\\Delta(E_{\\rm perp}) = 4 \\pi \\sigma$. There\nare two analogous boundary conditions here; one is \n\\begin{equation}\n\\label{eq:israel}\n\\Delta({da \\over\n du}) \\propto \\rho_{brane};\n\\end{equation}\ni.e. the change in the derivative of the scale factor as one crosses\nour 3-brane is given by the energy density on the brane.\n\nThe combination of the 5D Einstein equations with these boundary\nconditions leads to the modified Friedmann equations on our observable\nthree dimensional universe. Depending on what is in the bulk\n(off the brane), one can obtain a variety of possibilities such\nas $H^2 \\sim \\rho^n$ \\cite{cf}. Though we were at first concerned\nthat such modifications would be bad for cosmology, in fact\nthey can explain the observed acceleration of the universe today.\n\nCurrently one of our pressing goals would be to find a simple\nfive-dimensional energy momentum tensor that produces the Cardassian\ncosmology. While we succeeded in constructing an ugly toy model,\nit obviously does not represent our universe. Unfortunately,\ngoing from four to five dimensions is not unique, and is difficult.\n\n\\section{Motivation for Cardassian cosmology: 2) Fluid Description}\n\nAnother pressing goal is to find testable predictions of Cardassian\ncosmology. We want to be certain that ordinary astrophysical\nresults are not affected in an adverse way; e.g., we want energy\nmomentum to be conserved. In the interest of addressing these\ntypes of questions, we developed a four dimensional fluid description\nof Cardassian cosmology \\cite{gondolo}. This may be an effective\npicture of higher dimensional physics, or may be completely\nindependent of it. This fluid description may serve as a second\npossible motivation for Cardassian cosmology.\n\nHere we use the ordinary Friedmann equation, $H^2 = 8 \\pi G \\rho\/3$\nand the ordinary four dimensional Einstein equation,\n\\begin{equation}\nG_{\\mu\\nu} = 8 \\pi G T_{\\mu\\nu} .\n\\end{equation}\nWe take the energy density to be the sum of two terms:\n\\begin{equation}\n\\rho_{tot} = \\rho_M + \\rho_k\n\\end{equation}\nwhere $\\rho_k$ is a Cardassian contribution.\nAccompanying the two terms in the energy density are two\npressure terms:\n\\begin{equation}\np_{tot} = p_M + p_k ,\n\\end{equation}\nwhere the thermodynamics of an adiabatically expanding universe tell us that\n\\begin{equation}\np_k = \\rho_M \\left({\\partial \\rho_k \\over \\partial \\rho_M} \\right)_S\n- \\rho_k .\n\\end{equation}\nThese total energy density and pressure terms are what enter\ninto the four dimensional energy momentum tensor:\n\\begin{equation}\nT_{\\mu\\nu} = {\\rm diag}(\\rho_{tot},p_{tot},p_{tot},p_{tot}).\n\\end{equation}\n\nEnergy conservation is then automatically guaranteed:\n\\begin{equation}\nT_{\\mu\\nu;\\nu} = 0 .\n\\end{equation}\nThis equation implies a modified continuity equation,\nand a modified Euler's equation. We also require particle number\nconservation. Poisson's equation is also modified and becomes,\nin the Newtonian limit,\n\\begin{equation}\n\\nabla^2 \\phi = 4 \\pi G (\\rho_{tot} + 3 p_{tot}) .\n\\end{equation}\n\nAs an example, we can show the fluid description of the \npower law Cardassian model. We have\n\\begin{equation}\n\\rho_k = b \\rho_M^n \n\\end{equation}\nand\n\\begin{equation}\np_k = - (1-n) \\rho_k ,\n\\end{equation}\nwhich is a negative pressure for $n<2\/3$. The Poisson equation\nbecomes\n\\begin{equation}\n\\nabla^2 \\phi = - 4 \\pi G \\left[\\rho_M - (2-3n) b \\rho_M^n \\right]\n\\end{equation}\nOne can show that this simplified model runs into trouble\non galactic scales.\nThere is a new force ${d\\vec{v} \\over dt}|_{\\rm new}\n= - {\\vec{\\nabla} p_k \\over \\rho_M}$ that\ndestroys flat rotation curves. The fluid power law\nCardassian must be thought of as an effective model which applies\nonly on large scales. Hence, we proposed another possible\nmodel, the Modified Polytropic Cardassian model described earlier.\nDue to the existence of the additional parameter $q$, which\nis important on small scales, rotation curves are fine.\n\n\\subsection{Phantom Cardassian}\n\nSome Cardassian models can satisfy the dominant energy condition $w =\np_{tot}\/\\rho_{tot} > -1$ even with a dark energy component $w_k =\np_k\/\\rho_k < -1$, since both the ordinary and Cardassian components of\nthe energy density are made of the same matter and radiation. For\nexample, the MP Cardassian model wtih n=0.2 and q=2 has $w_k < -1$.\nThere is some evidence that $w_k < -1$ matches the data \\cite{mmot}.\n\n\\subsection{Speculation: Self Interacting Dark Matter?}\n\nWe here speculate on an origin for the new Cardassian term\nin the total energy density in the fluid model. The dark matter\nmay be subject to a new, long-range confining force (fifth force)\n\\begin{equation}\nF(r) \\propto r^{\\alpha -1},\\,\\,\\, \\alpha>0 .\n\\end{equation}\nThis may be analagous to quark confinement that exhibits negative\npressure. \n\nOur basic point of view is that Cardassian cosmology gives\nan attractive phenomenology, in that it is very efficient:\nmatter can provide all the observed behavior of the universe.\nTherefore it is worthwhile to examine all the avenues we\ncan think of as to the origin of these modified Friedmann terms\n(or alternatively of the fluid model).\n\n \n\\section{Discussion}\n\nWe have presented $H^2 = g(\\rho)$ as a modification to the Friedmann\nequation in order to suggest an explanation of the recent acceleration\nof the universe. In the Cardassian model, the universe can be flat\nand yet matter dominated. We have found that the new Cardassian\nmodifications can dominate the expansion of the universe after\n$z_{car} = \\mathcal{O}$$(1)$ and can drive an acceleration. We have\nfound that matter alone can be responsible for this behavior. The\ncurrent value of the energy density of the universe is then smaller\nthan in the standard model and yet is at the critical value for a flat\ngeometry.\nSuch a modified Friedmann equation may result from the existence of\nextra dimensions. Further work is required to find a simple\nfundamental theory responsible for Eq.(\\ref{eq:new}). A second\npossible motivation for Cardassian cosmology, namely a fluid\ndescription, has been developed \\cite{gondolo}, which allows us to\ncompute everything astrophysical. Comparison with SN observations,\nboth current and upcoming, were made and shown in Figure 1.\nIn future work, we plan to complete\na study of density perturbations in order to allow us to make\npredictions with observations, e.g. of cluster abundance or the ISW\neffect so that we may separate this model from others in upcoming\nobservations.\n\n\\section*{Acknowledgments}\n\nThis paper reflects work with collaborators Matt Lewis, Paolo Gondolo,\nYun Wang, Josh Frieman, Chris Savage, and Nori Sugiyama. We thank Ted\nBaltz for reminding us to consider the effects of the mass of the tau\nneutrino. I acknowledge support from the Department of Energy via the\nUniversity of Michigan and the Kavli Institute for Theoretical Physics\nat Santa Barbara.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nIn several recent papers (Lai \\& Tsang 2009; Tsang \\& Lai 2009c; Fu \\&\nLai 2011; Horak \\& Lai 2013), we have presented detailed study of the\nlinear instability of non-axisymmetric inertial-acoustic modes [also\ncalled p-modes; see Kato (2001) and Wagoner (2008) for review] trapped \nin the inner-most region of black-hole (BH) accretion discs. This global instability\narises because of wave absorption at the corotation resonance (where\nthe wave pattern rotation frequency matches the background disc\nrotation rate) and requires that the disc vortensity has a positive\ngradient at the corotation radius (see Narayan et al.~1987, Tsang \\& Lai 2008 \nand references therein). The disc vortensity (vorticity divided by\nsurface density) is given by\n\\be\n\\zeta={\\kappa^2\\over 2\\Omega\\Sigma},\n\\label{eq:zeta}\\ee \nwhere $\\Omega(r)$ is the disc rotation frequency, $\\kappa(r)$ is the \nradial epicyclic frequency and $\\Sigma(r)$ is the surface density\n\\footnote{Equation (\\ref{eq:zeta}) applies to barotropic discs in \nNewtonian (or pseudo-Newtonina) theory. See Tsang \\& Lai (2009c)\nfor non-barotropic discs and Horak \\& Lai (2013) for full\ngeneral relativistic expression.}. \nGeneral relativistic (GR) effect\nplays an important role in the instability: For a Newtonian disc, with\n$\\Omega=\\kappa\\propto r^{-3\/2}$ and relatively flat $\\Sigma(r)$\nprofile, we have $d\\zeta\/dr<0$, so the corotational wave absorption\nleads to mode damping. By contrast, $\\kappa$ is non-monotonic near a\nBH (e.g., for a Schwarzschild BH, $\\kappa$ reaches a maximum at\n$r=8GM\/c^2$ and goes to zero at $r_{\\rm ISCO}=6GM\/c^2$), the\nvortensity is also non-monotonic. Thus, p-modes with frequencies such\nthat $d\\zeta\/dr>0$ at the corotation resonance are overstable. Our\ncalculations based on several disc models and Paczynski-Witta \npseudo-Newtonian potential (Lai \\& Tsang 2009, Tsang \\& Lai 2009c)\nand full GR (Horak \\& Lai 2013) showed that the lowest-order\np-modes with $m=2,3,4,\\cdots$ have the largest growth rates, with the\nmode frequencies $\\omega \\simeq \\beta m\\Omega_{\\rm ISCO}$ (thus giving\ncommensurate frequency ratio $2:3:4,\\cdots$), where the dimensionless\nconstant $\\beta\\lo 1$ depends weakly the disc properties.\nThese overstable p-modes could potentially explain the \nHigh-frequency Quasi-Periodic Oscillations (HFPQOs) observed in BH\nX-ray binaries (e.g., Remillard \\& McClintock 2006; Belloni et al.~2012).\n\nThe effects of magnetic fields on the oscillation modes of BH\naccretion discs have been investigated by Fu \\& Lai (2009, 2011, 2012)\nand Yu \\& Lai (2013). Fu \\& Lai (2009) showed that the the basic wave\nproperties (e.g., propagation diagram) of p-modes are not strongly\naffected by disc magnetic fields, and it is likely that these p-modes\nare robust in the presence of disc turbulence (see Arras, Blaes \\&\nTurner 2006; Reynolds \\& Miller 2009). By contrast, other diskoseismic\nmodes with vertical structure (such as g-modes and c-modes) may be\neasily ``destroyed'' by the magnetic field (Fu \\& Lai 2009) or\nsuffer damping due to corotation resonance (Kato 2003; \nLi et al.~2003; Tsang \\& Lai 2009a). Although a modest\ntoroidal disc magnetic field tends to reduce the growth rate of the\np-mode (Fu \\& Lai 2011), a large-scale poloidal field can enhance the instability \n(Yu \\& Lai 2013; see Tagger \\& Pallet 1999; Tagger \\&\nVarniere 2006). The p-modes are also influenced by the magnetosphere\nthat may exist inside the disc inner edge (Fu \\& Lai 2012).\n\nSo far our published works are based on linear analysis. While these\nare useful for identifying the key physics and issues, the nonlinear\nevolution and saturation of the mode growth can only be studied by\nnumerical simulations. It is known that fluid perturbations near the\ncorotation resonane are particularly prone to become nonlinear (e.g.,\nBalmforth \\& Korycansky 2001; Ogilvie \\& Lubow 2003). \nMoreover, real accretion discs are more complex than any\nsemi-analytic models considered in our previous works. Numerical MHD\nsimulations (including GRMHD) are playing an increasingly important\nrole in unraveling the nature of BH accretion flows (e.g., De Villiers\n\\& Hawley 2003; Machida \\& Matsumoto 2003; Fragile et al.~2007; Noble\net al.~2009, 2011; Reynolds \\& Miller 2009; Beckwith et al.~2008, 2009;\nMoscibrodzka et al.~2009; Penna et al.~2010; Kulkarni et\nal.~2011; Hawley et al.~2011; O'Neill et al.~2011; Dolence et al.~2012;\nMcKinney et al.~2012; Henisey et al.~2012). Despite much\nprogress, global GRMHD simulations still lag far behind observations,\nand so far they have not revealed clear signs of HFQPOs that are directly \ncomparable with the observations of BH X-ray binaries\n\\footnote{Henisey et al.~(2009, 2012) found evidence of\nexcitation of wave modes in simulations of tilted BH accretion disks.\nHydrodynamic simulations using $\\alpha$-viscosity (Chan 2009; O'Neill et al.~2009)\nshowed wave generation in the inner disk region by viscous instability\n(Kato 2001). The MHD simulations by O'Neill et al.~(2011) revealed\npossible LFQPOs due to disk dynamo cycles. \nDolence et al.~(2012) reported transient QPOs in the numerical models\nof radiatively inefficent flows for Sgr A$^\\star$. McKinney et\nal.~(2012) found QPO signatures associated with the interface between\nthe disc inflow and the bulging jet magnetosphere (see Fu \\& Lai\n2012). Note that the observed HFQPOs are much weaker than LFQPOs,\ntherefore much more difficult to obtain by brute-force simulations.}.\nIf the corotation instability and its magnetic counterparts studied in\nour recent papers play a role in HFQPOs, the length-scale involved\nwould be small and a proper treatment of flow boundary conditions is\nimportant. It is necessary to carry out ``controlled'' numerical\nexperiments to capture and evaluate these subtle effects.\n\nIn this paper, we use two-dimensional hydrodynamic simulations to\ninvestigate the nonlinear evolution of corotational instability of\np-modes. Our 2D model has obvious limitations. For example it does not\ninclude disc magnetic field and turbulence. However, we emphasize that\nsince the p-modes we are studying are 2D density waves with no\nvertical structure, their basic radial ``shapes'' and real frequencies\nmay be qualitatively unaffected by the turbulence (see Arras et\nal. 2006; Reynolds \\& Miller 2009; Fu \\& Lai 2009). Indeed, several\nlocal simulations have indicated that density waves can propagate in\nthe presence of MRI turbulence (Gardiner \\& Stone 2005; Fromang et\nal. 2007; Heinemann \\& Papaloizou 2009). Our goal here is to\ninvestigate the saturation of overstable p-modes and the their\nnonlinear behaviours.\n\n\n\\section{Numerical Setup}\n\nOur accretion disc is assumed to be inviscid and geometrically thin so that the\nhydrdynamical equations can be reduced to two-dimension with vertically\nintegrated quantities. We adopt an isothermal equation of state \nthroughout this study, i.e. $P=c_s^2\\Sigma$ where $P$ is the vertically integrated\npressure, $\\Sigma$ is the surface density and $c_s$ is the constant\nsound speed. Self-gravity and magnetic fields are neglected.\n\nWe use the Paczynski-Witta Pseudo-Newtonian potential (Paczynski \\&\nWitta 1980) to mimic the GR effect:\n\\be\n\\Phi=-\\frac{GM}{r-r_{\\rm S}},\n\\ee\nwhere $r_{\\rm S}=2GM\/c^2$ is the Schwarzschild radius. \nThe corresponding Keplerian rotation frequency and radial epicyclic frequency are\n\\ba\n&&\\Omega_{\\rm K}=\\sqrt{\\frac{GM}{r}}\\frac{1}{r-r_{\\rm S}},\\\\\n&&\\kappa=\\Omega_{\\rm K}\\sqrt{\\frac{r-3r_{\\rm S}}{r-r_{\\rm S}}}.\n\\label{eq:grkappa}\n\\ea\nIn our compuation, we will adopt the units such that \nthe inner disc radius (at the Inner-most Stable Circular Orbit or ISCO)\nis at $r=1.0$ and the Keperian frequency at the ISCO is\n$\\Omega_{\\rm ISCO}=1$. In these units, $r_{\\rm s}=1\/3$, and\n\\be\n\\Omega_{\\rm K}=\\frac{2}{3}\\frac{1}{r-r_{\\rm s}}\\frac{1}{\\sqrt{r}}.\n\\ee\nOur computation domain extends from $r=1.0$ to $r=4.0$ in the radial\ndirection and from $\\phi=0$ to $\\phi=2\\pi$ in the azimuthal\ndirection. We also use the Keplerian orbital period\n($T=2\\pi\/\\Omega_{\\rm K}=2\\pi$) at $r=1$ as the unit for time.\nThe equilibrium state of the disk is axisymmetric. The surface density\nprofile has a simple power-law form\n\\be\n\\Sigma_0=r^{-1},\n\\ee\nwhich leads to a positive vortensity gradient in the inner disc region.\nThe equilibrium rotation frequency of the disc is given by \n\\be\n\\Omega_0(r)=\\sqrt{\\frac{4\/9}{r(r-1\/3)^2}-\\frac{c_s^2}{r^2}}.\n\\ee\nThroughout our simulation, we will adopt $c_s=0.1$ so that\n$\\Omega_0 \\simeq \\Omega_{\\rm K}$.\n\n\\begin{figure*}\n\\begin{center}\n$\n\\begin{array}{ccc}\n\\includegraphics[width=0.33\\textwidth]{m2vr1.ps} &\n\\includegraphics[width=0.33\\textwidth]{m3vr1.ps} &\n\\includegraphics[width=0.33\\textwidth]{mrvr1.ps}\n\\end{array}\n$\n\\caption{Evolution of the radial velocity amplitude $|u_r|_{\\rm max}$ \n(evaluated at $r=1.1$) for three runs with initial azimuthal mode number $m=2$\n(left panel), $m=3$ (middle panel) and \nwith random perturbations (right panel). \nThe dashed lines are the fits for the\nexponential growth stage (between $\\sim 10$ and $\\sim 30$ orbits) of\nthe mode amplitude.}\n\\label{fig:growth}\n\\end{center}\n\\end{figure*}\n\n\n\\begin{figure*}\n\\begin{center}\n$\n\\begin{array}{cc}\n\\includegraphics[width=0.45\\textwidth]{m2svr.ps} & \n\\includegraphics[width=0.45\\textwidth]{m3svr.ps} \\\\\n\\includegraphics[width=0.45\\textwidth]{m2sdvphi.ps} &\n\\includegraphics[width=0.45\\textwidth]{m3sdvphi.ps} \n\\end{array}\n$\n\\caption{Comparison of the radial profiles of velocity perturbations\n from non-linear simulation and linear mode calculation. The top and\n bottom panels show the radial and azimuthal velocity perturbations,\n respectively. The left and right panels are for cases with\n azimuthal mode number $m=2$ and $m=3$, respectively. In each panel,\n the dashed line is taken from the real part of the complex\n wavefunction obtained in linear mode calculation, while the solid line is from\n the non-linear simulation during the exponential growth stage (at\n $T=20$ orbits), with the quantities evaluated at $\\phi=0.4\\pi$.\n Note that the normalization factor is given in the $y$-axis label\n for the nonlinear simulation results. }\n\\label{fig:vrvphi}\n\\end{center}\n\\end{figure*}\n\nWe solve the Euler equations that govern the dynamics of the disc flow\nwith the PLUTO code \\footnote{publicly available at\n http:\/\/plutocode.ph.unito.it\/} (Mignone et al. 2007), which is a\nGodunov-type code with multiphysics and multialgorithm modules. For\nthis study, we choose a Runge-Kutta scheme (for time integration) and\npiecewise linear reconstruction (for space integration) to achieve\nsecond order accuracy, and Roe solver as the solution of Riemann\nproblems. The gird resolution we adopt is $(N_r\\times\nN_{\\phi})=(1024\\times 2048)$ so that each grid cell is almost\na square. Our runs typically last for 100 orbits (Keplerian orbits at\ninner disc boundary). To compare with the linear mode calculations\n(Lai \\& Tsang 2009), the inner disk boundary is set to be reflective\nwith zero radial velocity. At the outer disc boundary, we adopt the\nnon-reflective boundary condition. This is realized by employing the\nwave damping method (de Val-Borro et al. 2006) to reduce wave reflection. \nThis implementation mimics the outgoing radiative boundary condition \nused in the linear analysis.\n\n\n\\section{Results}\n\nWe carried out simulations with two different types of initial\nconditions for the surface density perturbation. In the first, we\nchoose $\\delta \\Sigma(r, \\phi)$ to be randomly distributed in $r$ and\n$\\phi$. In the second, we impose $\\delta \\Sigma(r, \\phi)\\propto\n\\cos(m\\phi)$ that is randomly distributed in $r$, so that the perturbation\nhas an azimuthal number $m$.\nIn all cases, the initial surface density perturbation has a small amplitude \n($|\\delta\\Sigma\/\\Sigma_{0}| \\leq 10^{-4}$).\n\n\\begin{table*}\n\\caption{Comparison of results from linear and nonlinear studies of overstable disc p-modes}\n\\centering\n\\begin{threeparttable}\n\\begin{tabular}{ c c c c c c c c }\n\\noalign{\\hrule height 1pt}\n$m$\\tnote{a} & $\\omega_{r}$\\tnote{b} &$\\omega_i$\\tnote{c} & $\\omega_{r}\/m\\Omega$\\tnote{d} & $\\omega_{r1}$\\tnote{e} & $|\\omega_r-\\omega_{r1}|\/\\omega_r$\\tnote{f} & $\\omega_{r2}$\\tnote{g} & $|\\omega_{r2}-\\omega_{r1}|\/\\omega_{r2}$\\tnote{h} \\\\\n\\hline\n$2$ & $1.4066$ & $0.0632$ & $0.7033$ & $1.3998$ & $0.5\\%$ & $1.4296$ & $2.1\\%$ \\\\\n$3$ & $2.1942$ & $0.0733$ & $0.7314$ & $2.1997$ & $0.3\\%$ & $2.2445$ & $2.0\\%$ \\\\\n$4$ & $3.0051$ & $0.0763$ & $0.7512$ & $2.9996$ & $0.2\\%$ & $3.0594$ & $2.0\\%$ \\\\\n$5$ & $3.8294$ & $0.0751$ & $0.7659$ & $3.7995$ & $0.8\\%$ & $3.8886$ & $2.3\\%$ \\\\\n$6$ & $4.6621$ & $0.0714$ & $0.7770$ & $4.6494$ & $0.3\\%$ & $4.7749$ & $2.6\\%$\\\\\n$7$ & $5.5007$ & $0.0664$ & $0.7858$ & $5.4992$ & $0.03\\%$ & $5.6756$ & $3.1\\%$ \\\\\n$8$ & $6.3436$ & $0.0607$ & $0.7930$ & $6.3492$ & $0.09\\%$ & $6.3189$ & $0.5\\%$\\\\\n\\noalign{\\hrule height 1.2pt}\n\\end{tabular}\n\\begin{tablenotes}\n\\item[a] Azimuthal mode number\n\\item[b] Mode frequency from the linear calculation (in units of Keplerian\n orbital frequency at the inner disc boundary; same for $\\omega_i$,\n $\\omega_{r1}$ and $\\omega_{r2}$)\n\\item[c] Mode growth rate from the linear calculation \n\\item[d] Ratio of wave pattern speed to the Keplerian orbital frequency at the inner disc boundary\n\\item[e] Mode frequency during the exponential growth stage of nonlinear simulation \n(peak frequency of the power density spectrum between $\\sim 10$ orbits and $\\sim 30$ orbits)\n\\item[f] Difference between $\\omega_r$ (linear result) and $\\omega_{r1}$ (nonlinear result)\n\\item[g] Mode frequency during the saturation stage of nonlinear simulation \n(peak frequency of the power density spectrum between $\\sim 30$ orbits and $\\sim 100$ orbits)\n\\item[h] Difference between $\\omega_{r1}$ and $\\omega_{r2}$\n\\end{tablenotes}\n\\end{threeparttable}\n\\label{tab:tab1}\n\\end{table*}\n\nFig.~\\ref{fig:growth} shows the evolution of the radial velocity\namplitude near the inner disc radius for runs with initial azimuthal\nmode number $m=2$, $m=3$ and random initial perturbation,\nrespectively. This velocity amplitude is obtained by searching for the\nmaximum $|u_r|$ at a given $r$ by varying $\\phi$ for each given time\npoint. We chose $r=1.1$ because this is where the largest radial\nvelocity perturbation is located (see the upper panels of\nFig.~\\ref{fig:vrvphi}). We see that in all cases there are three\nstages in the amplitude evolution. The first stage occupies roughly\nthe first $10$ orbits, during which the initial perturbation starts to\naffect the flow and presumably excites many modes\/oscillations in the\ndisc. In the second stage (from $\\sim 10$ orbits to $\\sim\n30$ orbits), the fastest growing mode becomes dominant and undergoes \nexponential growth with its amplitude increased by about four orders\nof magnitude. In the last stage (beyond $\\sim 30$ orbits), the\nperturbation growth saturates and its amplitude remains at\napproximately the same level. A fit to the exponential growth stage\ngives the growth rate of the fastest growing perturbation, which is\nthe slope of the fitted straight line. For the $m=2$ run, we find \n$0.17\/\\log_{10}(e)\/2\\pi\\simeq 0.0637$ (in units of the orbital\nfrequency at $r=1$) as the growth rate, which is quite consistent with\nthe result from our linear eigenmode calculation $\\omega_{i}\\simeq\n0.0632$ (the imaginary part of the eigenfrequency; Lai \\& Tsang 2009;\nFu \\& Lai 2011; see Table~\\ref{tab:tab1}). For the $m=3$ run,\nour simulation gives $0.074$ as the mode growth rate, close to\n$\\omega_{i}=0.0733$ from the linear calculation. In\nFig.~\\ref{fig:vrvphi}, we plot the radial profile of the velocity\nperturbation at a given time during the exponential growth stage\nof the simulation, and we compare it with the wavefunctions obtained\nfrom the linear mode calculation. Note that in each panel of\nFig.~\\ref{fig:vrvphi} we have normalized the different sets of data so\nthat they have the same scale. The normalization factor has also been\nincluded in figure labels for easy recovery of the absolute value\nin the simuation. We can see that wavefunctions obtained from two studies \nagree quite well. The only obvious differences are near the outer disc boundary,\nwhich can be attributed to the fact that outer boundary conditions\nemployed in the numerical simulation and linear mode calculation are not\nexactly the same. Nevertheless, this agreement and the agreement in \nthe mode growth rate in Fig.~\\ref{fig:growth} confirm that our non-linear simulations \nindeed capture the same unstable disc p-modes as those revealed in linear\nperturbation analysis.\n\n\n\n\\begin{figure*}\n\\begin{center}\n$\n\\begin{array}{cc}\n\\includegraphics[width=0.4\\textwidth]{m3fft2.ps} &\n\\includegraphics[width=0.4\\textwidth]{mrfft2.ps} \\\\\n\\includegraphics[width=0.4\\textwidth]{m3fft3.ps} &\n\\includegraphics[width=0.4\\textwidth]{mrfft3.ps} \n\\end{array}\n$\n\\caption{Power density spectra of the radial velocity perturbations\nnear the disk inner boundary ($r=1.1$). Each panel shows the normalized FFT\n magnitude as a function of frequency. The left and right\n columns are for runs with initial $m=3$ and with random perturbation,\n respectively. In the top and bottom panels, the Fourier transforms are\n sampled for time periods of [10, 30] orbits and [30,\n 100] orbits, respectively.}\n\\label{fig:fft}\n\\end{center}\n\\end{figure*}\n\n\n\\begin{figure*}\n\\begin{center}\n$\n\\begin{array}{cc}\n\\includegraphics[width=0.4\\textwidth]{m2cvr.ps} & \n\\includegraphics[width=0.4\\textwidth]{m3cvr.ps} \\\\\n\\includegraphics[width=0.4\\textwidth]{m2cdvphi.ps} &\n\\includegraphics[width=0.4\\textwidth]{m3cdvphi.ps} \n\\end{array}\n$\n\\caption{Evolution of the radial profiles of velocity perturbations from\n the non-linear simulations. The top and bottom panels show the radial\n and azimuthal components of velocity perturbation, respectively. \n The left and right panels are for cases with azimuthal mode number\n $m=2$ and $m=3$, respectively. In each panel, the data are taken\n from points with fixed $\\phi=0.4\\pi$. Different line types represent\n different times during the simulation.}\n\\label{fig:cvrvphi}\n\\end{center}\n\\end{figure*}\n\n\\begin{figure*}\n\\begin{center}\n$\n\\begin{array}{ccc}\n\\includegraphics[width=0.33\\textwidth]{m2vr_turn0020.ps} &\n\\includegraphics[width=0.33\\textwidth]{m3vr_turn0020.ps} &\n\\includegraphics[width=0.33\\textwidth]{mrvr_turn0020.ps} \\\\\n\\includegraphics[width=0.33\\textwidth]{m2vr_turn0030.ps} &\n\\includegraphics[width=0.33\\textwidth]{m3vr_turn0030.ps} &\n\\includegraphics[width=0.33\\textwidth]{mrvr_turn0030.ps} \\\\\n\\includegraphics[width=0.33\\textwidth]{m2vr_turn0050.ps} &\n\\includegraphics[width=0.33\\textwidth]{m3vr_turn0050.ps} &\n\\includegraphics[width=0.33\\textwidth]{mrvr_turn0050.ps} \n\\end{array}\n$\n\\caption{Evolution of the radial velocity for runs with initial $m=2$\n (left), $m=3$ (middle) and random (right) perturbations,\n respectively. From top to bottom, the times are $T=20$, $30$ and $50$ orbits, \n respectively. Note that the color scale varies from panel to panel.}\n\\label{fig:vrcolor}\n\\end{center}\n\\end{figure*}\n\n\nTo explore the time variability of the flow, we carry out\nFourier transform of the radial velocity $u_r(r, \\phi, t)$ at fixed\n$r$ and $\\phi$ during different evolutionary stages. In\nFig.~\\ref{fig:fft} we show some examples of the resulting power\ndensity spectra (normalized to the maximum value of unity). Different\nrows correspond to different total sampling times and different\ncolumns represent runs with different initial surface density\nperturbation. In the left columns, the disk has an initial\nperturbation with azimuthal mode number $m=3$. A mixture of various\nmodes\/oscillations are excited in the flow. After about $10$ orbits of\nevolution, one of them (the fastest growing mode) has its oscillation\namplitude grown by a large amount such that it dominates over other\nmodes. This corresponds to the primary spike in the top-left\npanel. The other spikes in the same panel are harmonics of this\nprimary spike (see the labels of the dashed vertical\nlines). Table~\\ref{tab:tab1} shows that the frequency\\footnote{All the\n frequencies in this paper are angular frequencies unless otherwise\n noted.} of this fastest growing mode ($\\omega_{r1}$) differs\nfrom the frequency obtained in linear mode calculation by only\n$0.3\\%$, which again demonstrates the consistency of these two studies. \nAfter the perturbation saturates (bottom-left panel), we see\nthat the basic structure of the power density spectrum does not change\nmuch except that these spikes are not as ``clean'' as in the linear regime;\nthis is probably due to the interaction of different modes. Compared with the upper\npanel, the location of the primary spike ($\\omega_{r2}$ in\nTable~\\ref{tab:tab1}) is increased by $2.0\\%$. In\nTable~\\ref{tab:tab1} we also include the results from both linear mode\ncalculation and numerical simulation for modes with other mode number\n$m$. The comparison illustrates two main points: First, the\nfrequencies of the fastest growing modes during the exponentially\ngrowth stage of numerical simulations are exceptionally close to the\nlinear calculation results (differ by less than $1\\%$); second, the\nfrequencies of the fastest growing modes during the saturation stage\nare only slightly higher (except for the $m=8$ mode which shows lower\nfrequency) than ones during the exponential growing stage. These indicate\nthat the mode frequencies obtained in linear mode calculation are\nfairly robust and can be reliably applied in the interpretation of HFQPOs.\n\nIn the right columns of Fig.~\\ref{fig:fft}, the simulation starts with\na random initial perturbation which excites modes with various\n$m$'s. During the exponential growth stage, six modes stand out. By\nexamining the location of the corresponding spikes, we know that these\nare the $m=3$, $4$, $5$, $6$, $7$, $8$ modes. Although the\n$m=3$ mode seems to be the most prominent one in this figure, we note\nthat this is because for this particular run the initial random\nperturbation happens to contain more $m=3$ wave components. If we were\nto start the run with a different initial random perturbation, then\nthe relative strengths of those peaks would also be different (not\nnecessarily with $m=3$ being dominant).\n\nFig.~\\ref{fig:cvrvphi} shows the comparison of the velocity perturbations \nduring the linear stage ($T=20$ orbits), at the end of the linear\nstage ($T=30$ orbits) and during the saturation stage ($T=50$ orbits)\nfor simulations with different initial perturbations. At $T=20$\norbits, the oscillation mainly comes from the single fastest growing\nmode and has a smooth radial profile. At $T=30$ orbits, the perturbation\nstarts to saturate, and the oscillation now consists of many different modes,\nand its radial profile exhibits sharp variations at several locations.\nThese sharp features remain after the saturation ($T=50$ orbits).\n\nTo see this evolution from a different perspective, in\nFig.~\\ref{fig:vrcolor} we show the color contours of radial velocity\nfor runs with different initial perturbations (different columns) at\ndifferent times (different rows). We can clearly see that spiral waves\ngradually develop due to the instability. As the system evolves, sharp\nfeatures in velcoities emerge. At the end (the bottom \nrow) the spiral arms becomes more irregular, which is related\nto the emergence and interaction of multiple modes after saturation.\n\nThe sharp velocity variations shown in\nFigs.~\\ref{fig:cvrvphi}-\\ref{fig:vrcolor} suggest shock-like features.\nNote that since our simulations adopt isothermal equation of state\n(with $P\/\\Sigma=$constant), there is no entropy generation and shock\nformation in the strict sense. The radial velocity jump (see\nFig.~\\ref{fig:cvrvphi}) is comparable but always smaller than the\nsound speed. Nevertheless, these sharp features imply that\nwave steepening plays an important role in the mode saturation. \n\n\\section{Conclusions}\n\nWe have carried out high-resolution, two-dimensional hydrodynamical\nsimulations of overtstable inertial-acoustic modes (p-modes) in BH\naccretion discs with various initial conditions. The evolution of\ndisc p-modes exhibits two stages. During the first (linear) stage, the\noscillation amplitude grows exponentially. In the cases with a\nspecific azimuthal mode number ($m$), the mode frequency, growth rate\nand wavefunctions agree well with those obtained in our previous\nlinear mode analysis. In the cases with random initial perturbation,\nthe disc power-density spectrum exhibits several prominent\nfrequencies, consistent with those of the fastest growing linear modes\n(with various $m$'s). These comparisons with the linear theory\nconfirm the physics of corotational instability that drives disc\np-modes presented in our previous studies (Lai \\& Tsang 2009; Tsang \\&\nLai 2009c; Fu \\& Lai 2011, 2012). In the second stage, the mode growth\nsaturates and the disc oscillation amplitude remains at roughly a\nconstant level. In general, we find that the primary disc oscillation\nfrequency (in the cases with specific initial $m$) is larger than the\nlinear mode frequency by less than $4\\%$, indicating the robustness of\ndisc oscillation frequency in the non-linear regime. Based on the\nsharp, shock-like features of fluid velocity profiles, we suggest that\nthe nonlinear saturation of disc oscillations is caused by wave\nsteepening and mode-mode interactions.\n\nAs noted in Section 1, our 2D hydrodynamical simulations presented in\nthis paper do not capture various complexities (e.g., magnetic field,\nturbulence, radiation) associated with real BH accretion discs.\nNevertheless, they demonstrate that under appropriate conditions, disc\np-modes can grow to nonlinear amplitudes with well-defined frequencies\nthat are similar to the linear mode frequencies. A number of issues\nmust be addressed before we can apply our theory to the interpretation\nof HFQPOs. First, magnetic fields may play an important role in the\ndisc oscillations. Indeed, we have shown in previous linear\ncalculations (Fu \\& Lai 2011,2012) that the growth of p-modes can be\nsignificantly affected by disc toroidal magnetic fields. A strong,\nlarge-scale poloidal field can also change the linear mode frequency\n(Yu \\& Lai 2013). Whether or not these remain true in the non-linear regime is\ncurrently unclear. Second, understanding the nature of the inner disc boundary\nis crucial. Our calculations rely on the assumption that the inner disc edge\nis reflective to incoming spiral waves. In the standard disc model with zero-torque\ninner boundary condition, the radial inflow velocity is not negligible near the ISCO, \nand the flow goes through a transonic point. While the steep density and velocity\ngradients at the ISCO give rise to partial wave reflection (Lai \\& Tsang 2009),\nsuch radial inflow can lead to significant mode damping such that the net growth rates of\np-modes become negative. This may explain the absence of HFQPOs in the thermal\nstate of BH x-ray binaries. However, it is possible that the inner-most region of \nBH accretion discs accumulates significant magnetic flux and forms a magnetosphere.\nThe disc-the magnetosphere boundary will be highly reflective, leading to the \ngrowth of disc oscillations (Fu \\& Lai 2012; see also Tsang \\& Lai 2009b).\nFinally, the effect of MRI-driven disc turbulence on the p-modes\nrequires further understanding. In particular, turbulent viscosity \nmay lead to mode growth or damping, depending on the magnitude and the \ndensity-dependence of the viscosity (R. Miranda \\& D. Lai 2013, in prep).\n\n\n\\section*{Acknowledgements}\nThis work has been supported in part by the NSF grants AST-1008245,\nAST-1211061 and the NASA grant NNX12AF85G. WF also acknowledges the\nsupport from the Laboratory Directed Research and Development Program\nat LANL.\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzbwkp b/data_all_eng_slimpj/shuffled/split2/finalzzbwkp new file mode 100644 index 0000000000000000000000000000000000000000..873745692659765ee3df7cd0da370478c5484e4d --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzbwkp @@ -0,0 +1,5 @@ +{"text":"\\chapter*{Abstract}\n\n\nWe present a system for the classification of mountain\npanoramas from user-generated photographs followed by\nidentification and extraction of mountain peaks from those\npanoramas. We have developed an automatic technique that,\ngiven as input a geo-tagged photograph, estimates its FOV\n(Field Of View) and the direction of the camera using a\nmatching algorithm on the photograph edge maps and a rendered\nview of the mountain silhouettes that should be seen\nfrom the observer's point of view. The extraction algorithm\nthen identifies the mountain peaks present in the photograph\nand their profiles. We discuss possible applications in social\nfields such as photograph peak tagging on social portals,\naugmented reality on mobile devices when viewing a mountain\npanorama, and generation of collective intelligence systems\n(such as environmental models) from massive social\nmedia collections (e.g. snow water availability maps based\non mountain peak states extracted from photograph hosting\nservices).\n\\chapter*{Sommario}\n\n\nProponiamo un sistema per la classificazione di panorami montani e per l'identificazione delle vette presenti in fotografie scattate dagli utenti.\nAbbiamo sviluppato una tecnica automatica che, data come input la foto e la sua geolocalizzazione, stima il FOV (Field of View o Angolo di Campo) e l'orientamento della fotocamera. Questo avviene tramite l'applicazione di un algoritmo di matching tra la mappa degli edge (bordi) della fotografia e alle silhouette delle montagne che dovrebbero essere visibili dall'osservatore a quelle coordinate.\nL'algoritmo di estrazione identifica poi i picchi delle montagne presenti nella fotografia e il loro profilo. \nVerranno discusse alcune possibile applicazioni in ambito sociale come ad esempio: \nl'identificazione e tagging (marcatura) delle fotografie sui social network, \nrealt\u00e0 aumentata su dispositivi mobile durante la visione di panorami montani e\nla generazione di sistemi di intelligenza collettiva (come modelli ambientali) dalle enormi collezioni multimediali dei social network (p.es. mappe della disponibilit\u00e0 di neve e acqua sulle vette delle montagne, estratte da servizi di condivisione di immagini).\n\\chapter*{Acknowledgements}\n\n\nForemost, I would like to thank Prof. Marco Tagliasacchi and Prof. Piero Fraternali for the opportunity to work with them on a such interesting and exciting project, and the continuous support during the preparation of this thesis.\n\nI would like to thank also:\n\\begin{itemize}\n\\item\nDr. Danny Chrastina (\\url{http:\/\/www.chrastina.net}) for helping with the preparation of this document.\n\\item\nDr. Ulrich Deuschle (\\url{http:\/\/www.udeuschle.de}) for his kind permission of using his mountain panorama generating web tool.\n\\item\nGregor Brdnik (\\url{http:\/\/www.digicamdb.com}) for providing the database of digital cameras and their sensor sizes.\n\\item\nMiroslav Sabo (\\url{http:\/\/www.mirosabo.com}) for the photograph of the Matterhorn used in this thesis.\n\\item\nAll my friends and colleagues (who are too many to be listed) who have helped and contributed to this work.\n\\end{itemize}\n\nLastly, but most importantly, I would like to thank my parents (Alexey and Irina) for making it possible for me to get here, and Valentina for constant support and motivation.\n\n\n\n\n\n\\chapter{Introduction}\n\\label{Introduction}\n\\thispagestyle{empty}\n\n\\vspace{0.5cm}\n\n\\noindent\nThe most suitable paradigm for representing this work is probably passive crowdsourcing, a field that is not trivial to define and that is not even a subcategory of crowdsourcing in the strict sense, as its name can suggest. A possible definition of the passive crowdsourcing discipline can be the union of crowdsourcing, data mining and collective intelligence (computer science fields that have much in common but are slightly different). These differences are often hard to notice due to the youth of these concepts and so the presence of a lot of confusion. For this reason the purpose of this chapter is to define these concepts unambiguously and to define the problem statement of this work.\n\n\\section{Human Computation}\nThe idea of the computer computation goal has always been that which Alan Turing expressed in 1950:\n\\begin{quotation}\n\n\\noindent{\\emph{``\nThe idea behind digital computers may be explained by saying that these machines are intended to carry out any operations which could be done by a human computer.\n''}\n Alan Turing \\cite{Turing:1995:CMI:216408.216410}\n}\n}\n\\end{quotation}\nThough current progress in computer science brings automated solutions of more and more complex problems, this idea of computer systems able to solve any problem that humans can solve is however far from reality. There are still a lot of tasks that cannot be performed by computers, and those that could be but are preferred to be computed by humans for quality, time, and cost reasons. These tasks lead to the field of human computation: a field that is hard to give a definition to, in fact several definitions of the term can be found: the most general, modern and suitable for our needs of which is probably that extracted from von Ahn's dissertation:\n\\begin{quotation}\n\n\\noindent{\\emph{``\n... a paradigm for utilizing human processing power to solve \nproblems that computers cannot yet solve.\n''}\n Luis von Ahn \\cite{VonAhn:2005:HC:1168246}\n}\n}\n\\end{quotation}\nHuman computation can be thought of as several approaches varying with the tasks involved, the type of persons involved in task completion, the incentive techniques used, and what type of effort the persons are required to make. It must be said that the classification of human computation is not an ordinary hierarchy with parents and children, but is instead a set of related concepts not necessary including one another. So a possible taxonomy of human computation (seen as a list of related ideas) has been produced by combining definitions given by Quinn et al. \\cite{Quinn:2011:HCS:1978942.1979148} and Fraternali et al. \\cite{Fraternali:2012:PHL:2263310.2263839} can be:\n\\begin{itemize}\n\\item\n\\emph{Crowdsourcing}:\nthis approach manages the distributed assignment of tasks to an open, undefined and generally large group of executors. The task to be performed by the executors is split into a large number of microtasks (by the work provider or the crowdsourcing system itself) and each microtask is assigned by the system to a work performer, who executes it (usually for a reward of a small amount of money). The crowdsourcing application (defined usually by two interfaces: for the work providers and the work performers) manages the work life cycle: performer assignment, time and price\nnegotiation, result submission and verification, and payment. In addition to the web interface, some platforms offer Application Programming Interfaces (APIs),whereby third parties can integrate the distributed work management functionality into their custom applications. Examples of crowdsourcing solutions are Amazon Mechanical Turk and Microtask.com \\cite{Fraternali:2012:PHL:2263310.2263839}.\n\\item\n\\emph{Games with a Purpose (GWAPs)}:\nthese are a sort of crowsourcing application but with a fundamental difference in user incentive technique: the process of resolving a task is implemented as a game with an enjoyable user experience. Instead of monetary earning, the user motivation in this approach is the gratification of the playing process. GWAPs, and more generally useful applications where the user solves perceptive or cognitive problems without knowing, address task such as adding descriptive tags and recognising objects in images and checking the output of Optical Character Recognition (OCR) for correctness. \\cite{Fraternali:2012:PHL:2263310.2263839}.\n\\item\n\\emph{Social Computing}:\na broad scope concept that includes applications and services that facilitate collective action and social interaction online with rich exchange of multimedia information and evolution of aggregate knowledge \\cite{parameswaran2007social}. Instead of crowdsourcing, the purpose is usually not to perform a task. The key distinction between human computation and social computing is that social computing facilitates relatively natural human behavior that happens to be mediated by technology, whereas participation in a human computation is directed primarily by the human computation system \\cite{Quinn:2011:HCS:1978942.1979148}.\n\\item\n\\emph{Collective Intelligence}: if seen as a process, the term can be defined as groups of individuals doing things collectively that seem intelligent \\cite{malone-harnessing}. If it is seen instead as the process result, means the knowledge of any kind that is generated (even non consciously and not in explicit form) by the collective intelligence process. Quinn et al. \\cite{Quinn:2011:HCS:1978942.1979148} classifies it as the superset of social computing and crowdsourcing, because both are defined in terms of social behavior. The key distinctions between collective intelligence and human computation are the same as with crowdsourcing, but with the additional distinction that collective intelligence applies only when the process depends on a group of participants. It is conceivable that there could be a human computation system with computations performed by a single worker in isolation. This is why part of human computation protrudes outside collective intelligence \\cite{Quinn:2011:HCS:1978942.1979148}.\n\\item\n\\emph{Data Mining}:\nthis can be defined broadly as the application of specific algorithms for extracting patterns from data \\cite{Fayyad96knowledgediscovery}. Speaking about human-created data the approach can be seen as extracting the knowledge from a certain result of a collective intelligence process. Creating this knowledge usually is not the goal of the persons that generate it, in fact often they are completely unaware of it (just think that almost everybody contributes to the knowledge of what are the most popular web sites just by visiting them: they open a web site because they need it, not to add a vote to its popularity). Though it is a very important concept in the field of collective intelligence, machine intelligence applied to social science and passive crowdsourcing (that will be defined in the next section) is a fully automated process by definition, so it is excluded from the area of human computation.\n\\item\n\\emph{Social Mobilization}:\nthis approach deals with social computation problems where the timing and the efficiency is crucial. Examples of this area are safety critical sectors like civil protection and disease control.\n\\item\n\\emph{Human Sensors}:\nexploiting the fact that the mobile devices tend to incorporate more and more sensors, this approach deals with a real-time collection of data (of various natures) treating persons with mobile devices as sensors for the data. Examples of these applications are earthquake and other natural disaster monitoring, traffic condition control and pollution monitoring.\n\\end{itemize}\n\n\\section{Passive Crowdsourcing}\nThe goal of any passive human computation system is to exploit the collective effort of a large group of people to retrieve the collective intelligence this effort generates or other implicit knowledge. In this study case the collective effort is taking geo-tagged photographs of mountains and publishing them on the Web, the intelligence deriving from this collection of photographs is the availability of the appearances of a mountain through time, the knowledge we want to extract having these visual appearances of mountains is the evolution of its environmental properties in time (which can be for example snow or grass presence at a certain altitude) or using some ground truth data even to predict these features where the use of physical sensors for those measurements is difficult or impossible (i.e. snow level prediction at high altitudes). Passive crowdsourcing is an approach that is not trivial to classify within the described taxonomy: it can be best classified in an area that includes: \\emph{crowdsourcing} for the fact of exploiting the effort of human computation, \\emph{collective intelligence} since its extraction is the primary goal of the approach and \\emph{data mining} as it refers to the procedure of extracting some results of human computation from the public Web data. Figure \\ref{fig:crowdsourcingTaxonomy} shows the taxonomy proposed by Quinn et al. \\cite{Quinn:2011:HCS:1978942.1979148} with this proposal of passive crowdsourcing collocation.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=0.85\\columnwidth]{.\/figures\/crowdsourcing_taxonomy}\n\t\\caption{Proposed taxonomy of human computation including passive crowdsourcing.}\n\t\\label{fig:crowdsourcingTaxonomy}\n\\end{figure}\n\nThe main advantage of passive crowdsourcing with respect to traditional crowdsourcing is the enormous availability of data to collect and very reduced costs of its retrieval, a significant disadvantage, on the other hand, is the form of the data: not always perfectly suitable for the goals and almost always needing to be processed and analyzed before being used. In other words, traditional crowdsourcing asks the user to shape the data in the input in a certain way, paying for this; passive crowdsourcing instead retrieves the result of past work, but shaped as it is, so processing it next. Just think about a study of the public opinion trends of a certain commercial product, two different ways to collecting these opinions can be used:\n\\begin{itemize}\n\\item\n\\emph{Crowdsourcing}: a certain amount of money is paid to each person who gives his proper opinion about the product being studied (approach also called crowdvoting). This method guarantees perfectly shaped collected data: the form of the questions and answers can be decided and modeled easily, but obviously the availability of the data to be collected will be limited by the number of the people ready to take that survey, and the costs of collecing a huge dataset of opinions will not surely be low.\n\\item\n\\emph{Passive Crowdsourcing}: the publicly available web data is full of opinions of the customers of a certain product (think about a customer that posts a photograph of that product, or comments a photograph uploaded by someone else, declaring their personal opinion about that product), the cost of retrieving these photographs and comments are almost null and the availability of this content is enormous. The big problem anyway is the shape of this data: given a photograph the algorithm must decide whether it is pertinent to the study or not (if the object in the photograph is the product being looked for) and to estimate if the opinion expressed by the user is positive or negative.\n\\end{itemize}\n\nIn the context of this work the huge amount of available data is fundamental, so the passive crowdsourcing approach is prefered, and the purpose of this work is exactly to deal with the problem of data shaping and analyzing.\n\nAnother significant advantage of passive crowdsourcing (and very important for this work) is the availability of the implicit information the user himself is unaware of: if a person is asked whether the peak of the Matterhorn (a mountain in the Italian Alps) was covered with snow or not in August of 2010 - he does not remember, but an appropriate image processing technique can extract this information from the photograph the user posted on his social network profile during his vacation.\n\n\\section{User Generated Content}\nThe amount of available user generated media content on the Web nowadays is reaching unprecedented mass: Facebook alone hosts 240 billion photographs and gets 300 million new ones every day \\cite{Doherty2010Lifelog}. This massive input allows collective intelligence applications to reach unprecedented results in innumerable domains, from smart city scenarios to land and environmental protection.\n\nA significant portion of this public dataset are geo-tagged photographs. The availability of geo-tags in photos results from two current trends: first, the widespread habit of using the smartphone as a photo camera; second, the increasing number of digital cameras with an integrated GPS locator and Wi-Fi module. The impact of these two factors is that more and more personal photographs are being published on social portals and more and more of them are precisely geo-tagged \\cite{npdPhotosSmartphones}. A time-stamped and geo-located photo can be regarded as the state of one or more objects at a certain time. From a collection of photographs of the same object at different moments, one can build a model of the object, study its behavior and evolution in time, and build predictive models of properties of interest.\n\nHowever, a problem in the implementation of collective intelligence applications from visual user generated content (UGC) is the identification of the objects of interest: each object may change its shape, position and appearance in the UGC, making its tracking in a set of billions of available photographs an extremely difficult task. In this work we aim at realizing a starting point for the applications that generate collective intelligence models from user-generated media content based on object extraction and model construction in a specific environmental sector: the study of mountain conditions. We harness the\ncollective effort of people taking pictures of mountains from different positions and at different times of the year to produce models describing the states of selected mountains and their changing snow conditions over time. To achieve this objective, we need to address the object identification problem, which for mountains is more tractable than the general case described above, thanks to the fact that mountains are among the most motionless and immutable objects present on the planet. This problem of mountain identification in fact will be the goal of this work.\n\n\\section{Problem Statement}\nGiven a precisely geo-tagged photograph, the goal of this work is to determine whether the photograph contains a mountain panorama, if yes, estimate the direction of view of the photo camera during the shot and identify the visible mountain peaks on the photograph.\n\nWe will describe in the detail the proposed algorithm, how it has been implemented and tested with the result of successfully matching 64.2\\% of the input geo-tagged photographs.\n\nThe algorithm of peak detection that is presented can be used in applications where the purpose is to identify and tag mountains in photographs provided by users. The algorithm can also be used for the creation of mountain models and of their related properties, such as, for example, the presence of snow at a given altitude.\n\nTwo representative examples of the first type of usage are:\n\\begin{itemize}\n\\item Mountain peak tagging of user-uploaded photographs on photo sharing platforms that allow anyone browsing that photograph to view peak names of personal interest.\n\\item Augmented reality on mobile devices with real-time peak annotation on the device screen while in camera mode.\n\\end{itemize}\n\nAn example of the usage for environmental model building is the construction of a model for correcting ground and satellite based estimates of the Snow Water Equivalent (SWE) on mountains peaks (which is described in the Conclusions and Future Work section).\n\n\\section{Document Structure}\nIn the next chapter several past works will be discussed, each one relevant in its own way and field to this work: from passive crowdsourcing and influenza surveillance to image processing and vision-based UAV navigation.\n\nIn the third chapter the proposed algorithm itself is described in detail and an efficient vector-based matching algorithm is explained.\n\nIn chapter four the realized implementation of the discussed algorithm is explained, including all improvement techniques developed (even those that have been rejected after the validation phase).\n\nIn the fifth chapter the results of the tests performed on the implemented algorithm are listed with the data set structure and error metric description.\n\nFinally in the last chapter the conclusions about this work are drawn with the possible future direction of this project.\n\n\\chapter{Related Work and State of the Art}\n\\label{Related Work and State of the Art}\n\\thispagestyle{empty}\n\n\\vspace{0.5cm}\n\n\\noindent\nThis work combines many disciplines of social computing, image processing, machine intelligence and environmental modeling with the relative emblematic problems: passive crowdsourcing, object identification (in particular mountain boundaries) and pose estimation, collective intelligence extraction and knowledge modeling, snow level estimation. All of the these problems have been largely analyzed and studied recently, but only a very small part of them combines several of the listed problems.\n\nQuoting all the works in those fields would be impossible, so in this chapter several examples of works from each of these fields (often combined together) will be discussed, from social problems and those closely inherent to this work, to examples of non-social applications, emphasizing in fact the wide possibility of goals that can be reached with these approaches.\n\nThe fundamental concept of computational social science is treated by Lazer et at. \\cite{lazer2009life}, explaining the trend of social science, moving from the analysis of individuals to the study of society: nowadays almost any action performed, from checking email, making phone calls and checking our social network profile to going for a walk, driving to the office by car, booking a medical check-up or even paying at the supermarket with a credit card leaves a digital fingerprint. The enormous amount of these fingerprints generates data, that pulled together properly give an incredibly detailed picture of our lives and our societies. Understanding these trends of societal evolution and so also individuals changing is a complex task, but with an unprecedented amount of potential and latent knowledge waiting to be extracted. The authors discuss the problems of this approach such as managing of privacy issues (think about the NRC report on GIS data \\cite{national2007Putting} describing the possibility to extract the individual information even from well anonymized data, or the online health databases which were pulled down after a study revealed the possibility of confirming the identities \\cite{dnaBlocked}). The lack of the both approach and infrastructure standards in this emerging field is also discussed:\n\\begin{quotation}\n\n\\noindent{\\emph{``\nThe resources available in the social sciences are significantly smaller, and even the physical (and administrative) distance between social science departments and engineering or computer science departments tends to be greater than for the other sciences. The availability of easy-to-use programs and techniques would greatly magnify the presence of a computational social science. Just as mass-market CAD software revolutionized the engineering world decades ago, common computational social science analysis tools and the sharing of data will lead to significant advances. The development of these tools can, in part, piggyback on those developed in biology, physics and other fields, but also requires substantial investments in applications customized to social science needs.\n''}\n Lazer et at. \\cite{lazer2009life}\n}\n}\n\\end{quotation}\n\n\\section{Passive Crowdsourcing and Collective Intelligence}\nAn important work in this field (even if the authors do not use the term of passive crowdsourcing, but it is exactly the type of problem solving we mean this term for) is performed by Jin et al. \\cite{Jin:2010:WSM:1873951.1874196} in their study of society trends from photograph analysis. The authors propose a method for collecting information to identify current social trends, and also for the prediction of trends by analyzing the sharing patterns of uploaded and downloaded social multimedia. Each time an image or video is uploaded or viewed, it constitutes an implicit vote for (or against) the subject of the image. This vote carries along with it a rich set of associated data including time and (often) location information. By aggregating such votes across millions of Internet users, the authors reveal the wisdom that is embedded in social multimedia sites for social science applications such as politics, economics, and marketing \\cite{Jin:2010:WSM:1873951.1874196}.\n\nGiven a query, the relevant photographs with relative metadata are extracted from Flickr, and the global social trends are estimated. The motivation for introducing this approach is the low cost of crawling this information for the companies and the industries, as well as the possibility to analyze this data almost instantaneously with respect to common surveys. The implementation of the proposal is also discussed, with several tests in various fields that gave incredibly promising results:\n\\begin{itemize}\n\\item\n\\emph{Politics}: the popularity scores analysis of the candidates Obama and McCain during the USA president elections of 2008 gave the resulting trends which were correct within a tenth of a percent of the real election data.\n\\item\n\\emph{Economics}: a product distribution map (with \\emph{iPod} as the example) around the world over time was successfully drawn.\n\\item\n\\emph{Marketing}: the sales of past years of several products such as music players, computers and cellular phones were estimated and the results actually match the official sales trends of those products.\n\\end{itemize}\n\nAlthough being an innovative and efficient approach, it does not deal with the visual content of the images themselves (the authors also highlight this fact, declaring the intent to add this analysis in the future). An example of collective intelligence extraction using image content is the work of Cao et al. \\cite{conf\/icassp\/CaoLGJHH10} that presents a worldwide tourism recommendation system. It is based on a large-scale geotagged web photograph collection and aims to suggest with minimal input to tourists the destination they would enjoy. By taking more than one million of geotagged photographs, dividing them into clusters by geographical position and extracting the most representative photographs for each area, the system is able to propose destinations to the user having in input the set of keywords and images of the places the user likes. From a conceptual point of view this system tries to simulate a friend that knows your travel tastes and suggests destinations for your new journey, but with the difference that this virtual friend has visited millions of places all around the world. This lies exactly in the concept of collective intelligence: exploiting the effort of hundreds of thousands of people uploading their travel photographs, the authors build a knowledge on what the various world tourism destinations feel like.\n\n\\subsection{Real-Time Social Monitoring}\nAn approach related to the concept of Human Sensors treated in the previous chapter is the process of monitoring the online activity of the persons to predict in advance (or at least to identify quickly) the occurrence of some phenomena. The activity to be monitored and the phenomena to detect can be very different. The most popular online activity to monitor is for sure the web searches of the users, such as in works of Polgreen et al. \\cite{Polgreen_usinginternet}, Ginsberg et al. \\cite{citeulike:3681665} and Johnson et al. \\cite{15361003} (in Johnson et al. also the access logs to health websites are analyzed) that propose influenza surveillance exploiting web search query statistics, examining the connection between searches for influenza and actual influenza occurrence and finding strong relationships. Another example of web query monitoring is the prediction of macroeconomic statistics by Ettredge et al. \\cite{Ettredge:2005:UWS:1096000.1096010}, in particular the prediction of unemployment rate trend based on the frequency of search terms likely used by people seeking employment.\n\nAnother important source of activity are social networks, for example Twitter messages, used by Culotta \\cite{Culotta:2010:TDI:1964858.1964874} in the monitoring of influenza, similar to the references described above, and by Sakaki et al. \\cite{Sakaki:2010:EST:1772690.1772777} for detecting the earthquakes.\n\n\\section{Object Identification and Pose Estimation}\nThe pose estimation of the photograph will be the key problem of this work, even if the estimated variable is only the orientation and not the position (that is given in input), and is a very commonly treated problem. Several examples will be listed here with problems, each one with its own elements in common with this work.\n\nAn example of a relatively different problem of pose estimation with respect to this work is the estimation of the geographic position of a photograph proposed by Hays and Efros \\cite{Hays:2008:im2gps}: though the estimation of the position is performed with the analysis of the visual content of the image, a purely data-driven scene matching approach is applied to estimate a geographic area the photograph belongs to.\n\nRamalingam et al. \\cite{RBSB09} instead present a work that at first sight can seem the opposite of this one (instead of estimating the orientation given the position they estimate the position given the orientation: always perpendicular to the terrain) but it is very similar: the authors describe a method to accurately estimate the global position of a moving car using an omnidirectional camera and untextured 3D city models. The idea of the algorithm is the same: estimate the pose by matching the input image to a 3D model (city model in this case, elevation model in case of our work). The described algorithm extracts the skyline from an omni-directional photograph, generates the virtual fisheye skyline views and matches the photograph a to view, estimating in this way the position of the camera.\n\nOther works about pose estimation given 3D models of cities and buildings have been recently published, such as world wide pose estimation using 3D point clouds by Li et al. \\cite{Li:2012:WPE:2402940.2402943} in which the SIFT features are located and extracted in the photograph and matched with the worldwide models. The particular point of this pose estimation is that it does not use any geographical information, but it estimates the position, orientation and the focal length (and so the field of view).\n\nBaatz et al. \\cite{Baatz:2012:LCM:2125160.2125170} addresses the problem of place-of-interest recognition in urban scenarios exploiting 3D building information, giving in output the camera pose in real world coordinates ready for augmenting the cell phone image with virtual 3D information. Sattler et al. instead deals with the problem of the 2D-to-3D correspondence computation required for these cases of pose estimation of urban scenes, demonstrating that direct 2D-to-3D matching methods have a considerable potential for improving registration performance.\n\nAn innovative idea of photograph geolocalization by learning the relationship between ground level appearance and overhead appearance and land cover attributes from sparsely available geotagged ground-level images \\cite{LinCrossView} was introduced by Lin et al. \\cite{LinCrossView}: unlike traditional geolocalization techniques it allows the geographical position of an isolated photograph to be identified (with no other geotagged photographs available in the same region). The authors exploit two previously unused data sets: overhead appearance and land cover survey data. Ground and aerial images are represented using HoG \\cite{Dalal:2005:HOG:1068507.1069007}, self-similarity \\cite{shechtman2007matching}, gist \\cite{Oliva:2001:MSS:598425.598462} and color histograms features. For each of these data sets the relationship between ground level views and the photograph data is learned and the position is estimated by two proposed algorithms that are also compared with three other pre-existing techniques:\n\\begin{itemize}\n\\item\n\\emph{im2gps}:\nproposed by Hays and Efros \\cite{Hays:2008:im2gps} already described a few lines above, that does not make use of aerial and attribute information and can only geolocate query images in locations with ground-level training imagery.\n\\item\n\\emph{Direct Match (DM)}:\nmatches the same features for ground level images to aerial images with no translation, assuming that the ground level appearance and overhead appearance are correlated.\n\\item\n\\emph{Kernelized Canonical Correlation Analysis (KCCA)}:\nis a tool to learn the basis along the direction where features in different views are maximally correlated, used as a matching score. It however presents significant disadvantages: singular value decomposition for a non-sparse kernel matrix is need to solve the eigenvalue problem, making the process unfeasible as training data increases, and secondly, KCCA assumes one-to-one correspondence between two views (in contrast with the geolocalization problem where it is common to have multiple ground-level images taken at the same location) \\cite{Hays:2008:im2gps}. \n\\end{itemize}\n\nThe methods proposed by the authors are instead:\n\\begin{itemize}\n\\item\n\\emph{Data-driven Feature Averaging (AVG)}:\nbased on the idea that well matched ground-level photographs will tend to have also similar aerial and land cover attributes, this technique translates the ground level to aerial and attribute features by averaging the features of good scene matches.\n\\item\n\\emph{Discriminative Translation (DT)}:\nan approach that extends AVG with also a set of negative training samples, based on the intuition that the scenes with very different ground level appearance will have distinct overhead appearance and ground cover attributes (assumption obviously hypothetical and not always true).\n\\end{itemize}\n\nIn the performed tests the algorithm was able to correctly geolocate 17\\% of the isolated query images, compared to 0\\% for existing methods.\n\n\\subsection{Mountain Identification}\nAll the pose estimation works described so far in this chapter deal with urban 3D models, as does the majority of pose estimation research. This is not surprising since an accurate 3D model is fundamental for this kind of task, and the massive increase of the 3D data of buildings and cities in the last few years makes these studies possible. Apart from urban models however, there is another type of 3D data that is largely available, which evolved much earlier than urban data (even if usually with lower resolution): terrain elevation data. The elevation data, presented usually as a geographical grid with the altitude for each point, can be easily seen as a 3D model, and the most interesting objects formed by these models are for sure mountains. For this reason also mountains are sometimes the identified objects in the pose estimation task, here several examples of these works will be discussed.\n\nThe most significant work in this sector is probably that presented by Baboud et al. \\cite{Baboud2011Alignment}, which given an input geotagged photograph, introduces the matching algorithm for correct overlap identification between the photograph boundaries and those of the virtually generated panorama (based on elevation datasets) that should be seen by the observer placed in the geographical point where the photograph has been taken from. This algorithm, which will be discussed in detail further, is the starting point of this thesis work. It must be highlighted however, that the goal of the authors is peak identification for the implementation of photograph and video augmentation; this work instead aims to identify mountain peaks to extract the appearances of the mountain for environmental model generation.\n\nAnother important work, dealing not only with the direction, but also with the position estimation of a mountain image, was written by Naval et al. \\cite{Jr97estimatingcamera}. Their algorithm does not work with the complete edges of the image but only with the skyline extracted using a neural network. The position and the orientation is then computed by nonlinear least squares.\n\nA proposal that exploits the skyline and mountain alignment is the vision-based UAV navigation described by Woo et al. \\cite{WooSLKK07}, that with pose estimation in a mountain area (a problem similar to the other described works) introduces the possibility of vision-based UAV navigation in a mountain area using an IR (Infra-Red) camera by identifying the mountain peaks. This is a navigation method that is usually performed by extracting features such as buildings and roads, not always visible and available, so mountain peak and skyline recognition brings big advantages.\n\nA challenging task with excellent results is described by Baatz et al. \\cite{Baatz:2012:LSV:2403006.2403045} with a proposal of an algorithm that given a photograph, estimates its position and orientation on large scale elevation models (in case of the article a region of 40000~km$^2$ was used, but in theory the technique can be applied to position estimation on a world scale. The algorithm exploits the shape information across the skyline and searches for similarly shaped configurations in the large scale database. The main contributions are a novel method for robust contour\nencoding as well as two different voting schemes to solve the large scale camera pose recognition from contours. The first scheme operates only in descriptor space (it checks where in the model a panoramic skyline is most likely to contain the current query picture) while the second scheme is a combined vote in descriptor and rotation space \\cite{Baatz:2012:LSV:2403006.2403045}. The original six-dimensional search is simplified by the assumptions that the photograph has been taken close to the terrain level and the photograph has usually only a small roll with respect to the horizon (both assumptions were made also during the development of this work algorithm). Instead of supposing to have the right shot position, this technique renders the terrain view from a digital elevation model on a grid defined by distances of approximately 100~m~$\\times$~100~m for a total number of 3.5 million cubemaps (this is how the authors call the renders). The key problem of the work is in fact, the large scale search method to efficiently match the query skyline to one of the cubemaps. Given a query image, sky\/ground segmentation is performed following an approach based on unary data costs \\cite{Martin:2004:LDN:977249.977379,Luo:2002:PAD:839290.842704} for a pixel being assigned sky or ground. The horizon skyline then is represented by a collection of vector-quantized local contourlets (contour words, similar in spirit to visual words obtained from quantized image patch descriptors) that are matched to the collected cubemaps with a voting stage that retrieves the most probable cubemaps to contain the query skyline, the top 1000 candidates are then analyzed with geometric verification using iterative closest points to determine a full 3D rotation. Evaluation was performed on a data set of photographs with manually verified GPS tags or given location, and in the best implementation 88\\% of the photographs were localized correctly.\n\n\\section{Environmental Study}\nIn this section the problem of snow level and snow water equivalent (SWE) estimation will be described, with the current state of the art and used methods as well as the benefits that an environmental model based on passive crowdsourcing photograph analysis can bring.\n\nSnow Water Equivalent (SWE) is a common snowpack measurement. It is the amount of water contained within the snowpack. It can be thought of as the depth of water that would theoretically result if you melted the entire snowpack instantaneously. To determine snow depth from SWE you need to know the density of the snow. The density of new snow ranges from about 5\\% when the air temperature is 14\\ensuremath{^\\circ} F, to about 20\\% when the temperature is 32\\ensuremath{^\\circ} F. After the snow falls its density increases due to gravitational settling, wind packing, melting and recrystallization \\cite{whatIsSWE}. It is a very important parameter for industry and agriculture, and its correct estimation is one of the main problems facing the environmental agencies.\n\nThese measurements are usually performed with physical sensors and stations that are very sparse, so the need of a map of snow and SWE distribution introduces a key problem: the interpolation of this data. Interpolation is usually done (due to the lack of the sophisticated physical models dealing with altitude and temperature and the absence of other supporting data) with relatively rough methods such as:\n\\begin{itemize}\n\\item\nlinear or almost linear interpolation of the data, lacking of precision due to the low spatial density of the measurement stations and the physical laws and phenomena that change the snow level and density radically with changes in altitude\n\\item\ncombining the interpolated data with the satellite ground images, removing the snow estimation from the areas where the satellite image indicates its absence: although being a first step in image content analysis for the refining of snow data, it indicates only a binary presence or absence of the snow and has anyway a low spatial density due to its low resolution.\n\\end{itemize}\n\nThe availability of the estimation of snow properties retrieved from mountain photograph analysis, and the past years snow measure ground truth in the areas and altitudes different from those that are measured by physical sensors, can bring significant supporting data for a better interpolation and more precise snow cover maps.\n\n\\subsection{Passive Crowdsourcing and Environmental Modeling}\nAn important step in environmental modeling exploiting passive crowdsourcing was made by Zhang et al. \\cite{Zhang:2012:MPW:2187836.2187938} in a work that studies the problem of estimating geotemporal distributions of ecological phenomena using geo-tagged and time-stamped photographs from Flickr. The idea is to estimate a given ecological phenomena (the presence or absence of snow or vegetation in this case) on a given day at a given place, and generate the map of its distribution. A Bayesian probabilistic model was used, assuming that each photograph taken in the given time and place is an implicit vote to the presence or the absence of snow in that area. Exploiting machine learning, the probability that a photograph contain snow is estimated based both on the metadata of the photograph (tags) and the visual features (a simplified version of GIST classifier augmented with color features was used). The estimation of the results (which was made by the authors thanks to the fact that for the two phenomena studied, snowfall and vegetation cover, large-scale ground truth is available in the form of observations of satellites \\cite{Zhang:2012:MPW:2187836.2187938}) brought very promising results: a daily snow classification for a 2 year period for four major metropolitan areas (NYC, Boston, Chicago and Philadelphia) generates results with precision, accuracy, recall and f-measure equal to approximately $0.93$.\n\nEven if the algorithm proposed by authors uses a simple probabilistic model (presence or absence of snow and vegetation cover) it is an important introduction to the field of environmental and ecological phenomena analysis by mining photo-sharing sites for geo-temporal information about these phenomena. The goal of this work is to propose an image processing technique that brings these environmental studies to a new level.\n\\chapter{Proposed Approach}\n\\label{Proposed Approach}\n\\thispagestyle{empty}\n\n\\vspace{0.5cm}\n\n\\noindent\nIn this chapter we describe the proposed approach for the procedure of mountain identification, from the analysis of the properties of the photo camera used for the shot to the identification of the position in the photograph of each mountain peak.\n\nThe matching algorithm is partially based on an mountain contour matching technique proposed by Baboud et al. \\cite{Baboud2011Alignment}, so first this matching algorithm will be discussed with particular emphasis on its advantages and disadvantages and then our proposal will be described in detail.\n\n\\section{Analysis of Babound et al.'s algorithm}\nThe basic idea and the purpose of the algorithm is the same as that of this work: given a geotagged photograph in input, identify the mountain peaks present in it, by exploiting the elevation data and the dataset of peaks with their positions.\nThe algorithm in question treats the input image as spherical, and searches for the rotation on $SO(3)$ that gives the right alignment between the input and another spherical image representing the mountain silhouettes viewed by the shot point of the input photograph in all directions. This image is generated by the 3D terrain model based on a digital elevation map. A vector cross-correlation technique is proposed to deal with the matching problem: the edges of the images are modeled as imaginary numbers, the candidates for the estimated direction of view are extracted, and a robust matching algorithm is used to identify the correct match.\n\nAs will follow from the next sections, several changes with respect to the original algorithm have been introduced, most relevant among them are:\n\\begin{itemize}\n\\item\nInstead of dealing with the elevation models we rely on an external service, which generates the panorama view given the geographic coordinates and a huge set of parameters. The reason of this choice is the possibility to concentrate the work on matching technique, as well as the reliability and precision of a tool improved over the years.\n\\item\nInstead of supposing the field of view to be known, its estimation given a photograph and the basic EXIF information will be discussed together with the photograph scaling problem, that allows successful matching by a non scale-invariant algorithm.\n\\item\nInstead of searching for the orientation of the camera during the shot in three dimensions, we suppose that the observer's line of sight is always perpendicular to the horizon (no photograph in the test dataset had a significant tilt angle with respect to the horizon). This assumption simplifies the computational effort of the algorithm, by moving from considering the images as spherical to cylindrical.\n\\item\nOnce the camera direction is successfully estimated, the mountain peak alignment between the photograph and the panorama may still presents non-negligible errors due to inaccuracies of the estimated position of the observer. We will deal this problem, which is not treated in the original algorithm.\n\\end{itemize}\n\n\\section{Overview of proposed algorithm}\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=0.65\\columnwidth]{.\/figures\/algorithm_schema}\n\t\\caption{Schema of the mountain peak tagging algorithm.}\n\t\\label{fig:algorithmSchema}\n\\end{figure}\n\nThe process of analyzing a single photograph containing a mountain panorama to identify individual mountain peaks present in it consists of several steps, shown in Figure \\ref{fig:algorithmSchema}.\n\nFirst the geotagged photograph is processed in order to evaluate its Field Of View and to scale it making the photograph match precisely the rendered panorama image, next a 360-degree rendition of the panorama visible at the location is generated from the digital elevation models.\nAfter this the direction of the camera during the shot is estimated by extracting the edge maps of both photograph and rendered panorama images and matching them. Finally the single mountain peaks are tagged based on the camera angle estimation.\n\n\\section{Detailed description of the algorithm}\n\n\\subsection{Render Generation}\nThe key idea of the camera view direction estimations lies in generating a virtual panorama view of the mountains that should be seen by the observer from the point where the photograph was taken, and then matching the photograph to this panorama.\nThe generation of this virtual view is possible due to the availability of terrain elevation datasets that cover a certain geographical area with a sort of grid identifying the terrain elevation in each grid point. Depending on the source of the elevation data, the precision and spatial density of this grid can vary significantly, but usually it is very precise, reaching even a spatial density of 3 meters in public datasets (such as USGS, http:\/\/ned.usgs.gov).\n\nThough the accuracy of elevation models is crucial for our purposes, since an exact render is the basis for a correct match, we are not necessary looking at extremely high-resolution datasets as the mountains we are going to generate are usually located at a significant distance from the observer (starting from few hundreds of meters to tens of kilometers). For this project the use of an external service that generated these rendered panoramas was preferred instead of creating them from scratch using the elevation data. The advantage of this choice is the possibility to avoid dealing directly with digital elevation maps and with optical and geometric calculations, in order to provide the panorama exactly as it should be seen by a human eye or photo camera lens. Such tools thus free the system from laborious calculations, allowing simply to choose the parameters needed for the generation of the panorama such as the observer's position, altitude and angle of gaze.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/mountainview2}\n\t\\caption{Example of wrong position estimation problems with view positions and the relative generated panoramas. 1 - Original viewer's position and altitude, 2 - estimated position and wrong altitude, 3 - elevated estimated position used for the final panorama generation.}\n\t\\label{fig:mountainView}\n\\end{figure}\n\nThe choice of the observer's altitude deserves to be mentioned separately, as it is a critical point: clearly the ideal value that can be set (with an ideal elevation model) is the real altitude of the observer during the shot (which, after studying the available datasets, we consider more than legitimate to suppose to be on the terrain surface, so we estimate the altitude of the observer as the terrain altitude in that position). This information is readily available, but it brings some problems due to the uncertainty of the geotag of the photograph: let us imagine an observer standing on a peak or a ridge of a mountain - he has a broad view in front of him, but it is enough for him to take a few steps back and the panorama he was viewing becomes completely covered by the facade of the mountain in front of him. The same issue occurs with the photographs: a photograph taken from a peak or a ridge of a mountain (very frequently the case in public domain collections of mountain photographs). The error of estimating the photo camera position of few meters can lead to an overlay of a significant portion of the panorama. An intuitive technique that can walk around this problem is to add some constant positive offset to the estimated altitude to \"raise\" the observer above the terrain obstacles that have appeared due to the errors in position estimation. Figure \\ref{fig:mountainView} shows in a simplified way the problem and its resolution with the corresponding generated panoramas.\n\n\\subsection{Field of View Estimation and Scaling}\nUnless a scale invariant matching technique is going to be used (and it is not the case of this algorithm), once the input image and an image representing the expected panorama view are ready, the first problem is that in order to be correctly matched the objects contained in both images (in our case the objects to be matched are mountains) must have the same pixel size. Since we assume that both represent the view of an observer from the same geographical point, the same pixel dimension of the object involves also the same angular dimensions (the angle a certain object occupies in the observer's view). So we define our first photograph analysis problem as searching for the right scaling of the input photograph to have the same angular dimensions for the same mountains present in the photograph and the rendered panorama. We can write down the problem as finding a scale factor $k$ such that $$k\\frac{s_{p}}{a_{s}} = \\frac{s_{r}}{a_{r}}$$ where $s_{p}$ and $a_{p}$ are respectively the pixel size and the angular size of the input photograph and $s_{r}$ and $a_{r}$ are similarly the pixel size and the angular size of the rendered panorama.\n\nWe expect the panorama to be exact, so we consider the mountains to have the same width\/height ratio both on the photograph and on the panorama, so the relationship described above can be equivalently applied both to the horizontal and vertical dimensions: we will work with the horizontal dimenstion because the angular width of the panorama does not need any calculation (it is always equal to the round angle, $2\\pi$). Defining the Field Of View (FOV) of an image as its angular width we can rewrite the relationship as $$k*\\frac{w_{p}}{FOV_{s}} = \\frac{w_{r}}{FOV_{r}} = \\frac{w_{r}}{2\\pi}$$ where $w$ and $FOV$ stands for the pixel width and the FOV of respectively the photograph and the rendered panorama.\n\n\\begin{figure}[h!]\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/fov}\n\t\\caption{A simplified schema of digital photo camera functioning.}\n\t\\label{fig:fovSchema}\n\\end{figure}\n\nBefore explaining how to estimate the FOV of the input photograph we must introduce some brief concepts regarding digital photo camera structure and optics laws: in a very simplified way a photo camera can be seen as a lens defining the focus of the projection of a viewed prospect on a sensor that captures this projection. The size of the sensor is a physical constant and property of a photo camera, the so called focal length instead (that defines the distance between the sensor and the lens) usually varies with the changing of the optical zoom of the camera. The FOV of the captured image in this case is obviously the angular size of the part of the prospect projected on the sensor. Figure \\ref{fig:fovSchema} shows this simple schema. We can easily write the relationship between the FOV of the photograph and the properties of the photo camera at the instant of shooting ($s$ for sensor width and $l$ for focal length): $$FOV = 2\\arctan\\frac{s}{2l}$$\n\nCombining this definition with the previous relationship we can express the scaling factor $k$ that must be applied to the photograph in order to have the same object dimensions as: $$k = FOV\\frac{w_{r}}{2\\pi w_{p}} = \\frac{w_{r}}{\\pi w_{p}}\\arctan\\frac{s}{2l}$$\n\nScaling estimation is a purely mathematical procedure, and the quality of the results depends directly on the precision of the geotag and the accuracy of the rendered model (with exact GPS location and a render based on correct elevation data the scaling produces perfectly matchable images).\n\n\\subsection{Edge Detection}\n\nThe key problem of the algorithm is to perform matching between the photograph containing the mountains and the automatically generated drawing representing the same mountain boundaries: in Figure \\ref{fig:mountainOverlapping} we can see an example of a fragment of the photograph and a fragment of the rendered panorama. Both represent the same mountain seen from the same point of view, and in fact when we try to match them manually they overlap perfectly. In spite of this overlapping, the choice of the technique for their matching is not trivial. Even if the problem of matching (choosing the position of the photograph with respect to the panorama, maximizing the overlap between the mountains) will be discussed in the next sections, here we will briefly mention the possible techniques to evaluate the correctness of the overlap and the reasons for the necessity of the edge extraction procedure.\n\n\\begin{figure}[h!]\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/mountain_overlapping}\n\t\\caption{An example of a matching problem with the photograph fragment (top right), the panorama fragment (top left) and their overlapping (bottom).}\n\t\\label{fig:mountainOverlapping}\n\\end{figure}\n\nWe can see the matching problem as a classic image content retrieval problem with the photograph as the input image and the collection of the fragments of the rendered panorama (each one corresponding to a possible alignment position) as the set of available images to search in. The most similar image in the set will be considered as the best matching position and identified as the direction of view of the camera during the shot.\n\nSo what is the similarity measure that can be used in order to perform this image content retrieval problem? Clearly global descriptors (color and texture descriptors) \\cite{citeulike:10106398} are not the best choice: it is enough to look at the Figure \\ref{fig:mountainOverlapping} to understand that the panorama is always generated in gray-scale tones, with textures defined only by the terrain elevation, while the mountains in the photograph can be colored in different ways, and the textures are defined by a lot of details such as snow on the mountains and other foreground objects such as grass, stones and reflections on the water. Local image descriptors instead (for example, local feature descriptors such as SIFT \\cite{Lowe:2004:DIF:993451.996342} and SURF \\cite{Bay:2008:SRF:1370312.1370556}) seem slightly suitable for the needs of mountain matching (as discussed also by Valgren et al. in \\cite{Valgren:2010:SSS:1715935.1716080}, where the use of SIFT and SURF descriptors is highlighted for the matching of outdoor photographs in different season conditions), but even if they are good for matching the photographs containing the same objects in different color conditions, the matching between an object photograph and its schematic representation (in our case a rendered drawing) tends to fail. Even if very accurate, the model of a photograph will not generate local features with the same precision that another photograph of the same object would do: no local descriptors tested have been able to find the match between the two example images in Figure \\ref{fig:mountainOverlapping}.\n\n\\begin{figure}[h!]\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/mountain_edges_overlapping}\n\t\\caption{An example of a matching problem with the photograph edge fragment (top right), the panorama edge fragment (top left) and their overlapping (bottom).}\n\t\\label{fig:mountainEdgesOverlapping}\n\\end{figure}\n\n\nThe perfect overlap between the images in the example figure however brings up to the idea, that instead of traditional image descriptors we should match the boundaries of the images (this assumption has also an intuitive motivation: the contours are the most significant and time invariant properties of the mountains), so the next step in photograph and panorama processing is edge extraction from both the images. The result of edge map extraction from the images in Figure \\ref{fig:mountainOverlapping} is represented in Figure \\ref{fig:mountainEdgesOverlapping}. The matching procedure on the edge maps can now be seen as a cross-correlation problem \\cite{wiki:CrossCorrelation}.\n\nIn order to make the cross-correlation matching more sophisticated with a couple of techniques that will be presented later, the edge extraction component must accept as input an image (the photograph or the rendered panorama) and produce as output a 2D matrix of edge points, each one corresponding to the point on the input image with a strength (value between 0 and 1 representing the probability of the point to be an edge) and a direction (value between 0 and 2$\\pi$ representing the direction of the edge in that point).\n\n\\subsection{Edge Filtering}\nOnce we have generated the edge maps of the photograph we must deal with the problem of the noise edge points. A noise edge point can be defined as an extracted edge point representing a feature that is not present on the rendered panorama, in this case the edge point will not be useful for the match and will only be able to harm it. In other words a noise edge point is an edge point that does not belong to a mountain boundary (the only features represented on our rendered panoramas). This doesn't mean that any edge point belonging to a mountain is not a noise point, i.e. the edge points that define the snow border of a mountain are noise points because they do not belong to the mountain boundary but define a border that due to its nature can easily change from photograph to photograph and furthermore will not be present in the rendered view.\n\nHowever, edge points can be also present on the mountain surface; most of them usually belong to foreign objects. Mountains tend to be very edge-poor and detail-poor objects and often are placed in the background of a photograph with other edge-rich objects such as persons, animals, trees, or houses in the foreground. Let us think about a photograph of a person next to a tree with a mountain chain in the background: the mountains themselves will generate few edge points since they tend to be very homogeneously colored; foreground objects instead will generate a huge amount of edge points (which will be noise) because they are full of small details, each leaf of the tree and detail of the person's clothes will produce noise points.\n\nTaking into account the example of the edge extraction made in the previous chapter, we manually tag the extracted points as noise and non-noise points following the edges present in the panorama: the result is represented in the Figure \\ref{fig:filteringGoodBadBefore} and with a simple image analysis script we find that the the noise edge points are more than 90\\% of all extracted points.\n This value grows up to 98\\% when we deal with photographs with even more objects in foreground. Even if the matching algorithm we are going to propose includes penalizing the noise points in order to keep up with this problem the amount of the noise edges reaches such high levels that it cause almost any algorithm to fail simply from a statistical point of view (the intense density of noise points in some area tends to ``attract'' the matching position in that area), so an edge filtering technique is needed.\n \n \\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=0.8\\columnwidth]{.\/figures\/filtering_good_bad_before}\n\t\\caption{Noise (red) and non-noise (green) edge points on our example photograph.}\n\t\\label{fig:filteringGoodBadBefore}\n\\end{figure}\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=0.8\\columnwidth]{.\/figures\/filtering_good_bad_after}\n\t\\caption{Noise (red) and non-noise (green) edge points on our example photograph after the filtering procedure has been applied.}\n\t\\label{fig:filteringGoodBadAfter}\n\\end{figure}\n\nOne of the possible approaches is the detection of the skyline (the boundary between the sky and the objects present in the photograph) as implemented by Naval Jr et al. \\cite{Jr97estimatingcamera} for the analogue problem. Skyline detection is a non-trivial task that has been widely studied with several proposed approaches available in particular sectors and scenarios, such as skyline detection of a perspective infrared image by synthesizing a catadioptric image proposed by Bazin et al. \\cite{DBLP:conf\/icra\/BazinKDV09} or the extraction of the skyline from a moving omnidirectional camera as developed by Ramalingam et al. \\cite{RBSB09}. The skyline however is not the only boundary of our interest, the boundaries of mountains contained ``inside'' other boundaries are also important and significant for the matching, so we decided to opt for softer a filtering approach.\n\nWe have opted for a simple but efficient technique based on the intuitive assumption that the mountains are usually placed above the other objects on the photographs (with some exceptions such as clouds or other atmospheric phenomena): prioritizing (increasing the strength of) the higher points on the photograph with respect to the lower points. An example of the result of the edge filtering procedure is displayed in Figure \\ref{fig:filteringGoodBadAfter}. Most of the noise edges are filtered with almost all good edges intact (also the edges that would be cut with a skyline detection). The rate of the noise edge points dropped from 90\\% to 40\\% in this case.\n\n\\subsection{Edge matching (Vector Cross Correlation)}\nFirstly we define $C_{r}$ as the cylindrical image generated from the rendered panorama with the same height, and the base perimeter equal to the panorama's width, and $C_{p}$ as the input photograph projected on an empty cylindrical image of the same dimensions of $C_{r}$. Imagining two cylinders to be concentric the matching problem is defined as the two dimensional space search (the vertical position and the rotation angle of $C_{r}$ with respect to $C_{p}$), which leads to the best overlap between the mountains present on the cylindrical images. As the cylindrical images are only a projection of a rendered image on a cylinder, the problem can be equivalently seen as a search for the best overlap of two rectangular images defined by two integer numbers representing the coordinate offset of the photograph with respect to the panorama (obviously carrying out the horizontal overflow part of the photograph to the opposite side of the render, simulating a cylinder property) as shown in Figure \\ref{fig:cylinderMatching}.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/cylinder_matching}\n\t\\caption{Representation of cylindrical image matching and the equivalent rectangular image matching.}\n\t\\label{fig:cylinderMatching}\n\\end{figure}\n\nAs already introduced in previous sections, the matching will be performed with the edge maps of the images so we define the result of the edge detection algorithm as a 2D real-valued vector where each value is defined by $\\rho$ (the strength\/absolute value of the edge in the corresponding point of the input image) and $\\theta$ (the direction\/argument of that edge). Let $p(\\omega)$ and $r(\\omega)$ be the 2D real-valued vectors generated by edge detection of the photograph and the panorama render respectively. Then as proposed by Baboud et al. \\cite{Baboud2011Alignment} the best matching is defined as the position maximizing the likelihood between two images. This likelihood is defined as $$L(p,r) = \\int_{S^{2}} M(p(\\omega),r(\\omega))d\\omega$$\nwhere $M$ is the angular similarity operator:\n$$M(v_1,v_2) = \\rho ^{2}_{v_1} \\rho ^{2}_{v_2} \\cos 2(\\theta_{v_1} - \\theta_{v_2})$$\n\nThis technique of considering the edge maps as vector matrices and applying the cross-correlating resolution is called Vector Cross Correlation (VCC). The cosine factor is introduced in order to handle edge noise by penalizing differently oriented edges: the score contribution is maximum when the orientation is equal, null when they form a $\\frac{\\pi}{4}$ angle, and minimum negative when the edges are perpendicular (Figure \\ref{fig:penalizingEdges}). This penalty avoids that random noise edges contribute in a positive way to a wrong match position (a step in this direction was already made during the edge filtering).\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/penalizing_edges}\n\t\\caption{Example of an overlapping position with positive score with almost parallel edge intersection (left circle) and penalizing almost perpendicular edge intersection (right circle).}\n\t\\label{fig:penalizingEdges}\n\\end{figure}\n\nOne of the main advantages of this likelihood score is the possibility of applying the Fourier transform in order to reduce drastically the computation effort of best matching position: let us define $\\hat{p}$ and $\\hat{r}$ as 2D Fourier transforms of respectively $p(\\omega)$ and $r(\\omega)$. The VCC computation equation becomes\n\\begin{equation} \\label{eq:FFTVCC}\nL(p,r) = Re\\{\\hat{p}^{2} \\bar{\\hat{r}}^{2}\\}\n\\end{equation}\nThis cosine similarity VCC can be seen as the whole algorithm to perform the matching or only as its first step: if we build the 2D real-valued distribution of the likelihood values estimated in each possible position we can consider not only the best match but extract top-N peaks of the score distribution as the result candidates and then evaluate them with a more sophisticated and heavier technique in order to identify the correct match. One of these refining algorithms is the robust silhouette map matching metric technique designed by Baboud et al. \\cite{Baboud2011Alignment} which considers the edge maps as sets of singular connected edges. The likelihood of an overlap is the sum of the similarity between each edge $e_p$ of the photograph and the edge $e_r$ of the rendered panorama, where $e_r$ is enriched with a certain $\\epsilon_e$ neighborhood, $l$ represents the distance for which $e_p$ stays inside the $\\epsilon_e$ neighborhood of the $e_r$ and the similarity is defined as\n\\begin{equation} \\label{eq:robustMatching}\nM(e_p, e_r) = \n\\left\\{\\begin{matrix}\n0\n\\\\ \nl^a\n\\\\ \n-c\n\\end{matrix}\\right.\n\\begin{array}{ll}\n\\text{if }l = 0\n\\\\ \n\\text{if }(l > l_{fit}) \\land (e_p \\text{ enters and exits on the same side})\n\\\\ \n\\text{if }(l < l_{fit}) \\land (e_p \\text{ enters and exits on different sides})\n\\end{array}\n\\end{equation}\nwhere $a$, $c$ and $l_{fit}$ are predefined constants. The nonlinearity implied by the exponent $a$ makes longer edge overlaps receive more weight than the set of small overlaps, the constant $c$ instead introduces also in this metric the penalizing factor for the intersections of the edges considered wrong as in the VCC metric.\n\n\n\\subsection{Mountain Identification and Tagging}\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/result_wrong_peaks}\n\t\\caption{Example of a fragment of a matching result between the photograph (blue) and the panorama (red).}\n\t\\label{fig:resultWrongPeaks}\n\\end{figure}\n\nOnce the edges are matched and the direction of view of the camera during the shot is estimated, supposing to have the list of mountain peaks with their names and coordinates on the rendered panorama we can estimate the position of these peaks also on the photograph, even if it is not as trivial as can be initially thought. Intuition suggests that once the photograph is matched with the panorama the coordinates of the mountain peaks are the same on both (obviously with a fixed offset equal to the matching position in case of the photograph) but it is wrong. Looking at a fragment of the result of edge matching (the best estimated position by VCC) of a photograph in Figure \\ref{fig:resultWrongPeaks}: the overlapping between the photograph and panorama edges is almost perfect on the left part but tends to be always worse moving to the right, which means that the two edge maps are slightly different (probably due to the error in position and altitude estimation) and objectively the matching proposed by the VCC seems to be the best that can be obtained. This situation occurs very often, and in general we can say that the resulting edge matching, even when successful, presents small errors in singular mountain peaks due to the imperfections of the generated panorama model. We propose a method for precise mountain peak identification: for each mountain peak present in the panorama we extract an edge map pattern both from the photograph and the panorama centered in the coordinate of the peak on the panorama and weighted with a certain kernel function $f(d)$ for a fixed pattern radius $r$ and $d$ representing a point distance from the peak coordinate with the following properties:\n$$\n\\begin{array}{ll}\nf(0) = 1\n\\\\ \nf(x) = 0 \\text{ }\n\\\\ \nf(x_1) \\geq f(x_2) \\text{ }\n\\end{array}\n\\begin{array}{ll}\n\n\\\\ \n\\forall x \\geq r\n\\\\ \n\\forall x_1, x_2 \\, \\mid \\, x_1 \\leq x_2\n\\end{array}\n$$\nOnce the two patterns are extracted the VCC procedure is once again applied in order to match them, having only the edges of the mountain peak we are processing and a few surrounding edges, the matching position is no longer influenced by the other peaks and is refined to the exact position, an example of this procedure performed on the matching result introduced in Figure \\ref{fig:resultWrongPeaks} in order to identify one of the most right-placed peaks is displayed in Figure \\ref{fig:peakIdentification}.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/peak_identification}\n\t\\caption{Peak identification procedure on the previous matching example.}\n\t\\label{fig:peakIdentification}\n\\end{figure}\n\n\n\\chapter{Implementation Details}\n\\label{Implementation Details}\n\\thispagestyle{empty}\n\n\\vspace{0.5cm}\n\n\\noindent\nIn this chapter the implementation of the matching algorithm proposed in the previous chapter is presented, from the general description of the system architecture and programming languages involved to the value of each parameter and constant introduced in the algorithm presentation.\n\n\\section{Overview of the implementation architecture}\nThe whole system can be split into two macro areas:\n\\begin{itemize}\n\\item\n\\emph{Photograph analysis, panorama and mountain peak list generation}: implemented in PHP with web interface due to the simplicity of interfacing with the external panorama generation service and interface versatility. Given a geo-tagged photograph (or a photograph with explicitly specified the geographic point of the shot) the virtual panorama view is generated centered at the same origin point and the properties both of the photograph and the photo camera are retrieved in order to let the next step calculate the field of view and scale factor.\n\\item\n\\emph{Matching and mountain identification}: implemented in MATLAB for the high suitability with image processing techniques. Given a photograph with all necessary information, the panorama and the mountain list with the relative coordinates on the panorama, edge extraction and filtering is performed, the camera direction is estimated, and finally individual mountains are identified and tagged.\n\\end{itemize}\n\nThe choice of the operating parameters used in the final implementation was defined in the testing and validation phase, using a data set and a developed evaluation metric. Although the parameter values will be specified as soon as they are introduced, the detailed description and the reasons of discarding the other values will be presented in the next chapter.\n\n\\section{Detailed description of the implementation}\n\n\\subsection{Render Generation}\nAs already mentioned in the implementation an external panorama service has been used: the mountain view generator of Dr. Ulrich Deuschle \\footnote{\\url{http:\/\/www.udeuschle.de\/panoramen.html}} (Figure \\ref{fig:udeuschleScreenshot}). The service accepts as input several parameters, most important of which are the geographic position of the observer, his altitude and altitude offset, the view direction with the field of view, zoom factor, elevation exaggeration and the range of sight. Latitude and longitude are equal obviously to the geo-tag of the input photograph, and the horizontal extension (field of view) to the round angle, $2\\pi$.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/udeuschle_screenshot}\n\t\\caption{A screenshot of the mountain panorama generation web service interface.}\n\t\\label{fig:udeuschleScreenshot}\n\\end{figure}\n\nThe choice of the altitude (in fact of the altitude offset) as already anticipated in the previous chapter is a problematic point: a big offset can lead to vertical distortion of the panorama with respect to reality, a small offset on the other hand can lead more easily to partial occlusion of the panorama. The choice was made for $auto + 10$~m altitude setting, chosen by trial and error. Regarding the altitude choice an important option of the service named ``Look for summit point automatically'' must be highlighted: by turning on this option the panorama is generated from the highest point of the terrain within 200~m. The advantage\/disadvantage of using this option is the same as increasing the altitude offset: avoiding the occlusion of the panorama against distorting the generated view. Even if this option is very interesting we set it to off as the distortion of the panorama due to the observer's position changing and not only his altitude was too detrimental for the algorithm.\n\nFor the most realistic panorama the zoom factor and elevation exaggeration factors are both set to $1$. All the other parameters are left to their default values.\n\nThis service covers the the Alps, Pyrenees and Himalayas mountain range systems with two elevation datasets:\n\\begin{itemize}\n\\item\n\\emph{Alps} by Jonathan de Ferranti, with the spatial resolution of $1$~arcseconds\n\\item\n\\emph{SRTM} by CGIAR, with the spatial resolution of $3''$\n\\end{itemize}\n\nThe output is implemented with dynamic loading of the panorama, so is generated as a set of images to be aligned horizontally to form the complete output image and the one image containing the mountain peak names to be overlapped with the result image. Despite the unavailability of the service API, after the study of JavaScript and AJAX scripts of the web interface, a PHP function that simulates the functioning of the interface and collects all the parts of the result image was implemented.\n\nThe generation of the full panorama at the resolution of $20$~pixel\/degree and the peaks names takes on average 1 minute.\n\n\\subsection{Field of View Estimation and Scaling}\nTo estimate the field of view of the input photograph, as described in the previous chapter, the only information that must be known is the focal length of the photograph and the sensor width of the photo camera. The focal length is a parameter specified directly in the EXIF format of the image (tag = FocalLength, 37386 (920A.H)) \\cite{Technical2002Exchangeable}. The sensor width instead is not specified directly in the EXIF since it is a property of a photo camera, so the manufacturer and the model of the camera are extracted from the EXIF and then the sensor size is retrieved (manufacturer tag = Make, 271 (10F.H); model tag = Model, 272 (110.H)) \\cite{Technical2002Exchangeable}. \n\nClearly, the necessity of a database of camera sensor sizes emerges. This information is usually scattered on web sites of manufactures and technical references of the cameras, so it is not easy to collect them in one place. We used a database kindly provided by ``Digital Camera Database'' \\footnote{\\url{http:\/\/www.digicamdb.com}} web site, containing information about the sensor size of more than 3000 digital cameras. The main problem in using it remains in the fact that the manufacturer and model names of the same cameras are not always equal between that written in the EXIF and that provided by the sensor database, a couple of these mismatching examples are provided in Table \\ref{tab:exifMismatching}.\n\n\\begin{table}[!h]\n \\centering\n \\begin{tabularx}{\\textwidth}{|X|X|X|X|}\n \\hline\n \\multicolumn{2}{|c|}{\\tabhead{EXIF}} &\n \\multicolumn{2}{c|}{\\tabhead{Database}} \\\\\n \\hline\n \\tabhead{Manufacturer} &\n \\tabhead{Model} &\n \\tabhead{Manufacturer} &\n \\tabhead{Model} \\\\\n \\hline\n Canon &\n Canon PowerShot SX100 IS &\n Canon &\n PowerShot SX110 IS \\\\\n \\hline\n SONY &\n DSC-W530 &\n Sony &\n Cybershot DSC W530 \\\\\n \\hline\n NIKON &\n E5600 &\n Nikon &\n Coolpix 5600 \\\\\n \\hline\n OLYMPUS IMAGING CORP. &\n SP560UZ &\n Olympus &\n SP 560 UZ \\\\\n \\hline\n \\end{tabularx}\n \\caption{Examples of differences in the names of manufacturers and models between the EXIF specifications and the digital camera database.}\n \\label{tab:exifMismatching}\n\\end{table}\n\nTo find the correct photo camera in the database from the EXIF specifications name a text similarity score between the names is used and the most similar name is chosen. The text similarity is calculated by the \\emph{similar\\_text} PHP function proposed by Oliver \\cite{DBLP:books\/daglib\/0077674}, after several steps of preprocessing:\n\n\\begin{enumerate}\n\\item\nBoth the manufacturer and model names of both the EXIF and the database items are transformed to lower case.\n\\item\nIf the manufacturer name contains ``nikon'' the name is set to ``nikon''.\n\\item\nIf the manufacturer name contains ``olympus'' the name is set to ``olympus''.\n\\item\nIf the model name contains the manufacturer name it is cut off from the model name.\n\\item\nThe text similarity score is performed between the concatenation of the manufacturer and the model of both EXIF and database items.\n\\end{enumerate}\n\nOnce the focal length and the sensor size are retrieved, they are annotated within the image, the matching algorithm will later use them to calculate the field of view and scale factor.\n\n\\subsection{Edge Detection}\nFor the edge extraction the \\emph{compass} \\cite{Ruzon:2001:EJC:505471.505477} edge detector has been used. It returns exactly the output the matching algorithm needs (for each point the edge strength and direction) and has been chosen due to its ability of exploiting the whole color information contained in the image and not only the grayscale components as classical edge detectors do.\nThe \\emph{compass} detector deals well also with a significant problem of edge detectors: when the image presents a subjective boundary between two regions with pixels close in color (due to overlapping objects), most edge detectors compute a weighted average of the pixels on each side, and since the two values representing the average color of two regions are close, no edges are found \\cite{Ruzon:2001:EJC:505471.505477}. The chosen detector deals well with this problem, that is very likely in the context of mountain photographs and the boundaries between the snowy mountain and blue sky.\n\nThe edge detection procedure is applied both to the input photograph and the generated panorama, with the following parameters chosen by trial-and-error: standard deviation of Gaussian used to weight pixels $\\sigma = 1$, threshold of the minimum edge strength to be considered $\\tau = 0.3$.\n\n\\subsection{Edge Filtering}\nThe edge filtering approach used in the implementation treats the columns of the image separately: each column is split into segments separated by the zero strength edge points, and each segment is then split into segments of length $n$ points. The points of each $i$-th segment (starting from the top) are then multiplied by the factor of $b^{i-1}$.\n\nThe implementation uses $k = 2$ and $b = 0.7$. It allows a good filtering of noise objects at the bottom of the photograph, and meanwhile deals well with mountain edges placed below the clouds. Several examples of filtering are presented in Figure \\ref{fig:filteringExamples}:\n\\begin{itemize}\n\\item\n\\emph{First row}: a photograph with clouds and a rainbow completely removed from the edges map (thanks to the correct parameters of the edge extraction algorithm).\n\\item\n\\emph{Second row}: a photograph with strongly contrasting clouds, the edges of the mountains below the clouds are still present even if reduced in strength.\n\\item\n\\emph{Third row}: a photograph with different mountains, one in front of another with a reduced visibility, even if reduced in strength all the mountains are present on the filtered edge map and the noise edges of the terrain in the foreground are filtered.\n\\end{itemize}\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/filtering_examples}\n\t\\caption{Examples of edge filtering (one for each row): the original photograph (left), the extracted edges (center) and the filtered edges (right).}\n\t\\label{fig:filteringExamples}\n\\end{figure}\n\n\\subsection{Edge matching (Vector Cross Correlation)}\nWith the edge maps of the photograph and the panorama the matching of the edges is performed by fast Fourier transformation, with the Formula \\ref{eq:FFTVCC}. The MATLAB implementation of the VCC is the following:\n\n\\begin{lstlisting}[frame=single]\nfunction VCC = ComputeVCC(SP, DP, SR, DR)\n\n\ndimP = size(SP);\ndimR = size(SR);\nSR = [ zeros(dimP(1), dimR(2)) ; SR ; zeros(dimP(1), dimR(2)) ];\nDR = [ zeros(dimP(1), dimR(2)) ; DR ; zeros(dimP(1), dimR(2)) ];\n\nCOMP = complex( SP .* cos(DP) , SP .* sin(DP) );\nCOMR = complex( SR .* cos(DR) , SR .* sin(DR) );\n\nVCC = rot90(real(ifft2(conj(fft2(COMR.^2)) .* fft2(COMP.^2, size(COMR, 1), size(COMR, 2)))),2);\n\nVCC = VCC( dimP(1) + 1 : 2 * dimP(1) );\n\\end{lstlisting}\n\nThe resulting matrix of this function is the VCC score evaluation for each possible overlapping position, and the maximum value of the matrix identifies the best overlap. Instead of just peaking up the maximum value, several additional techniques have been implemented and tested:\n\n\\emph{Result smoothing}: the idea is that the result score peak will have high values also in the neighborhood of its position, and the noise score peaks instead do not: in other words the peak corresponding to the right position will have be broader with respect to noise score peaks. So smoothing the result score distribution penalizes the noise peaks: it will reduce more the height of the noise peaks with respect to the correct peak. In practice, after several tests, the incorrectness of this assumption emerges. Due to the nature of the VCC technique of penalizing the non parallel intersection of the edges, even if the score in the correct position is high, it is sufficient to move only few points out to reduce drastically the score, so the smoothing procedure was not efficient. Figure \\ref{fig:smoothing} shows an example of the score matrix (projected onto the $X$--$Z$ plane for readability) of the photograph of the Matterhorn used in the previous examples: the correct position represents already the maximum of the distribution, and not only it is intuitively visible that the correct position score peak is not smoother than the others, but it is also shown that more smoothing always leads to a smaller difference between correct and incorrect matches. \n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/smoothing}\n\t\\caption{Smoothing of the VCC result distribution with smoothing factor $k$. The red arrow identifies the correct matching.}\n\t\\label{fig:smoothing}\n\\end{figure}\n\n\\emph{Robust matching}: as proposed in the algorithm chapter, the top-$N$ peaks of the score matrix are extracted, each one is evaluated to find the best matching position. The evaluation is performed with Formula \\ref{eq:robustMatching}, that has been implemented in a simplified version in the following way:\n\\begin{enumerate}\n\\item\nThe information about the direction of the edge points is removed, all the edge points with strength bigger than $0$ are considered to have the strength equal to $1$.\n\\item\nOn the panorama edge map each point that is located less or equal than $r$ points from any edge point in the original edge map is set to $1$. In this way the $r$-neighborhood of the panorama edges is generated.\n\\item\nA simple intersection between the new panorama edge map and the photograph edge map is performed in the overlapping position that must be evaluated.\n\\item\nThe resulting intersection (composed only by $0$ and $1$ strength points) is divided into clusters where any point of the cluster is located less than $d$ from at least one point of the same cluster.\n\\item\nThe Formula is applied where the clusters are singular edges, and the edge length is the number of the points in the cluster.\n\\end{enumerate}\n\nAfter performing tests on the dataset also this approach has been rejected as increasing drastically the computational effort of the matching without introducing significant benefits to the results.\n\n\\emph{Different scale factors}: since the reduced times of VCC computation with the fast Fourier transform ($\\sim 1$~second for the VCC score matrix generation in our implementation) an approach to reduce the impact of the wrong photograph shot position or the incorrect field of view estimation (due for example to a wrong sensor peaking) is to perform the matching at several scale levels and not only the estimated scale. In the current implementation several scaling intervals with respect to the estimated level were evaluated, and the scale factor with the best VCC maximum value was picked. Obviously a larger scale factor will lead to a larger photograph edge map and so bigger VCC score, so when comparing different maxim values an inverse quadratic scale factor must be applied to each score (inversely proportional to the area added to the image due to the scale changing).\n\n\\subsection{Mountain Identification and Tagging}\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/filtering_function}\n\t\\caption{Peak extraction: edge map before pattern extraction (left), implemented kernel function with $r = 1$ (center), edge map after pattern extraction (right).}\n\t\\label{fig:filteringFunction}\n\\end{figure}\n\n\nThe kernel function chosen for the peak extraction for the identification is the triweight function, defined as:\n\n$$\nf(d) = \n\\left\\{\\begin{matrix}\n\\left(1-\\left(\\frac{d}{r}\\right)^{2}\\right)^{3} \\\\ \n0\n\\end{matrix}\\right.\n\\begin{matrix}\n\\text{ for }d \\leq r \\\\\n\\text{ for }d > r\n\\end{matrix}\n$$\n\nWith $r = 200$. The kernel function plot and its effect on the peak extraction are shown in Figure \\ref{fig:filteringFunction}.\n\n\\chapter{Experimental Study}\n\\label{Experimental Study}\n\\thispagestyle{empty}\n\n\\vspace{0.5cm}\n\nIn this chapter the photograph direction estimation quality is investigated. The precision of the estimation is computed varying the operating parameters of the implementation and the type of the photographs contained in the data set.\n\n\\section{Data sets}\nThe analysis is conducted on a set of 95 photographs of the Italian Alps, collected from several photographers with geo-tag set directly by a GPS component or manually by the photographer, so we consider them very precisely geo-tagged. The photographs vary in several aspects: several examples are shown in Figure \\ref{fig:dataset}, and the categories of interest are composed as described in Table \\ref{tab:datasetCategoriesComposition}.\n\\begin{table}[!h]\n \\centering\n \\begin{tabularx}{\\textwidth}{|X|c|c|}\n \\hline\n \\tabhead{Category} &\n \\tabhead{Option} &\n \\tabhead{Data set portion} \\\\\n \\hline\n \\multirow{2}{*}{Source} &\n Photo camera &\n 38 \\% \\\\\n &\n Cellular phone &\n 62 \\% \\\\\n \\hline\n \\multirow{2}{*}{Cloud presence} &\n None &\n 41 \\% \\\\\n &\n Minimal &\n 29 \\% \\\\\n &\n Massive &\n 23 \\% \\\\ \n &\n Overcast &\n 7 \\% \\\\ \n \\hline\n \\multirow{2}{*}{Skyline composition} &\n Mountains and terrain only &\n 87 \\% \\\\\n &\n Foreign objects &\n 13 \\% \\\\\n \\hline \n \\end{tabularx}\n \\caption{Data set categories of interest composition.}\n \\label{tab:datasetCategoriesComposition}\n\\end{table}\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.00\\columnwidth]{.\/figures\/dataset}\n\t\\caption{Several photographs from the collected data set.}\n\t\\label{fig:dataset}\n\\end{figure}\n\n\n\\section{Operating parameters}\nThe operating parameters used in the algorithm described in the previous sections, and to be evaluated, are listed in Table \\ref{tab:datasetOperatingParameters}. The value in bold defines the default value that gives the best evaluation results and is used in the final implementation proposal.\n\n\\begin{table}[!h]\n \\centering\n \\begin{tabularx}{\\textwidth}{|X|c|c|}\n \\hline\n \\tabhead{Full name} &\n \\tabhead{Parameter} &\n \\tabhead{Tested values} \\\\\n \\hline\n Photograph edge strength threshold &\n $\\rho_p$ &\n 0.1,0.2,\\textbf{0.3},0.4,0.5,0.6,0.7 \\\\\n \\hline\n Panorama edge strength threshold &\n $\\rho_r$ &\n 0.1,\\textbf{0.2},0.3,0.4,0.5 \\\\\n \\hline \n Photograph edge filtering base &\n $b_p$ &\n 0.5,0.6,\\textbf{0.7},0.8,0.9,1.0 \\\\\n \\hline \n Panorama edge filtering base &\n $b_r$ &\n 0.5,0.6,0.7,0.8,0.9,\\textbf{1.0} \\\\\n \\hline \n Photograph edge filtering max segment length &\n $l_p$ &\n 1,\\textbf{2},3,4,5 \\\\\n \\hline \n Panorama edge filtering max segment length &\n $l_r$ &\n 1,2,3,4,5 \\\\ \n \\hline \n Photograph scaling interval &\n $\\pm k \\%$ &\n \\textbf{0},1,2,5,10 \\\\ \n \\hline \n \\end{tabularx}\n \\caption{Operating parameters (defaults in bold).}\n \\label{tab:datasetOperatingParameters}\n\\end{table}\n\n\\section{Evaluation metric}\nEach photograph in the data set and the corresponding generated panorama are manually tagged by specifying on both one or more pairs of points corresponding to the same mountain peaks. Once the direction of view is estimated and the best overlap of the edges is found, the error of the direction estimation is defined as the average distance between the position of the peak on the panorama and the position of the peak on the photograph projected on the panorama with the current overlap position. The distance is then expressed as the angle, considering the panorama width as $2\\pi$ angle. The error therefore lies between 0 and just over $\\pi$ angle (the worst horizontal error is $\\pi$ in the case of opposite directions, but since the vertical error is also counted the total error can theoretically slightly exceed this value).\n\nGiven a photograph and a panorama rendered with the resolution of $q$ pixel\/degree, both tagged with $N$ peaks, let $x_{pi}$\/$y_{pi}$ and $x_{ri}$\/$y_{ri}$ be respectively the coordinates (expressed in pixels) of the $i$-th peak on the photograph and the panorama image, the estimation error is defined as:\n\n$$\ne = \\frac{\\sqrt{(\\sum_{i = 1}^{N} (x_{pi} - x_{ri}))^2 + (\\sum_{i = 1}^{N} (y_{pi} - y_{ri}))^2 }}{Nq}\n$$\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.00\\columnwidth]{.\/figures\/errorOffsetValidation}\n\t\\caption{Schematic example of edge matching validation with photograph edges (red) and panorama edges (blue), all the three peaks are validation points.}\n\t\\label{fig:errorOffsetValidation}\n\\end{figure}\n\nAn important aspect of the average computing must be highlighted: instead of the standard distance error averaging by taking into account the singular distances, this error metric calculates the average of the both image coordinates axis components, and then the distance with these components is computed. In practice, a positive error offset compensates a negative one. This is made to compensate the imperfections between the scale of the mountains on the photograph and panorama: if the overlap presents both positive and negative offsets between peak pairs it means that the photograph and the panorama are differently scaled, so the best overlap position is the one that brings the sum of the offset components to zero. An example of this reasoning is shown as a schematic case of edge matching in Figure \\ref{fig:errorOffsetValidation} the validation points include all three peaks of the photograph and panorama view, but the different scaling prevents perfect overlapping and the offsets are both positive and negative. Thus the best possible overlapping position is the one shown, that makes the offsets opposite and making the sum of the offsets converge to zero.\n\nTherefore even when the validation points cannot all be matched perfectly, the technique of not taking the absolute value of the offsets makes it anyway possible to find the best match and obtain an offset of zero (perfect score).\n\nFor readability reasons all the validation errors from this point will be specified in arc degrees ($\\ensuremath{^\\circ}$).\n\nOnce the matching error is calculated for every photograph present in the data set, the general score of the run of the algorithm is defined as the percentage of the portion of the data set containing all the photographs with matching error below a certain threshold. This threshold in other words defines the limit of the matching error for a photograph to be considered matched correctly, and in the evaluation of this work is set to $4\\ensuremath{^\\circ}$.\n\n\\section{Evaluation results}\nWith the all operating parameters set to their default values the algorithm has correctly estimated the orientation of \\textbf{64.2\\%} of the total photographs. In particular the distribution of the matching error through the photographs of the data is presented in Figure \\ref{fig:mainResultHist}: it is clear that the choice of the error threshold has a limited impact on the correct matching rate since almost all the correctly matched photographs have error between $0\\ensuremath{^\\circ}$ and $2\\ensuremath{^\\circ}$.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.00\\columnwidth]{.\/figures\/main_result_hist}\n\t\\caption{Histogram of the number of the photographs with a certain validation error (split into intervals of width $1$~\\ensuremath{^\\circ}.)}\n\t\\label{fig:mainResultHist}\n\\end{figure}\n\nThe results obtained with the described data set decomposed into categories as described in Table \\ref{tab:datasetCategoriesComposition} are summarized in Table \\ref{tab:datasetCategoriesResults}.\n\n\\begin{table}[!h]\n \\centering\n \\begin{tabular}{|c|c|c|}\n \\hline\n \\tabhead{Category} &\n \\tabhead{Option} &\n \\tabhead{Correct matches} \\\\\n \\hline\n \\multirow{2}{*}{Source} &\n Photo camera &\n 72.2 \\% \\\\\n &\n Cellular phone &\n 59.3 \\% \\\\\n \\hline\n \\multirow{2}{*}{Cloud presence} &\n None &\n 71.8 \\% \\\\\n &\n Minimal &\n 67.9 \\% \\\\\n &\n Massive &\n 45.5 \\% \\\\ \n &\n Overcast &\n 66.7 \\% \\\\ \n \\hline\n \\multirow{2}{*}{Skyline composition} &\n Mountains and terrain only &\n 65.1 \\% \\\\\n &\n Foreign objects &\n 58.3 \\% \\\\\n \\hline \n \\end{tabular}\n \\caption{Results of the algorithm with respect to the data set categories of interest.}\n \\label{tab:datasetCategoriesResults}\n\\end{table}\n\nFrom the results of data set categories it follows that the photographs shot with a photo camera are more frequently aligned than those shot with a cellular phone: an explanation for this trend can be the fact that the photo camera photographers usually use the optical zoom that is correctly annotated in the focal length of the photograph, cellular phones instead usually allow only digital zoom, which leaves unchanged the annotated focal length and so leads to an incorrect Field of View estimation and therefore an incorrect alignment.\n\nThe trend of the result as cloud conditions vary follows a logical sense (more clouds leads to a higher number of noise edges and to worse performance) and reveals that the presence of clouds is a significant reason for failure of the algorithm, a problem that will be treated in the next chapter. A good score of the overcast case can be explained by the fact that when the sky is completely covered by clouds there are no edges between the real sky and clouds, so there are no noise edges (it can be thought of as the sky painted by a darker cloud color).\n\nThe presence of foreign objects such as trees, buildings, and persons in the skyline obviously penalizes the alignment, but as can be seen from the results it is not a critical issue.\n\nIn the next chapter several techniques that should improve the results of mountain identification are presented and discussed.\n\n\\subsection{Operating parameter evaluation}\nWe now present and comment on the effect of the individual parameters on the overall results:\n\n\\textbf{\\emph{Photograph edge strength threshold ($\\rho_p$)}}: The choice of the strength threshold for the photograph edges presents two trends: higher threshold means fewer noise edges, but lower threshold means to include also the weaker edges belonging not to the skyline but to secondary terrain edges. The dependency of the result on $\\rho_p$ is displayed in Figure \\ref{fig:opRhoP}, the optimal balance between the two trends is reached with $\\rho_p = 0.3$.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.00\\columnwidth]{.\/figures\/op_rho_p}\n\t\\caption{Overall algorithm results varying the $\\rho_p$ operating parameter.}\n\t\\label{fig:opRhoP}\n\\end{figure}\n\n\\textbf{\\emph{Panorama edge strength threshold ($\\rho_r$)}}: The choice of the strength threshold for the panorama edges is very similar to that of the photograph, with a difference regarding small thresholds: since the panorama is rendered by a program it does not present weak edges, which means that the introduction of a small threshold has almost no effect on the matching efficiency. The dependency of the result on $\\rho_r$ is displayed in Figure \\ref{fig:opRhoR}: in fact all tested values smaller or equal to $0.3$ bring the same (and optimal) result. In the implementation the parameter is set to $\\rho_r = 0.2$.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.00\\columnwidth]{.\/figures\/op_rho_r}\n\t\\caption{Overall algorithm results varying the $\\rho_r$ operating parameter.}\n\t\\label{fig:opRhoR}\n\\end{figure}\n\n\\textbf{\\emph{Photograph edge filtering base ($b_p$)}}: The choice of the filtering base for the photograph edges defines how strongly the lower edge points will be reduced in strength, $b_p = 1.0$ means that no filtering will be performed, and $b_p = 0.0$ that only the first segment of each row of the photograph will be left. The dependency of the result on $b_p$ is displayed in Figure \\ref{fig:opBP}: it is evident that filtering is absolutely needed since the result corresponding to $b_p = 1.0$ is significantly smaller than all the other values; the result stabilizes for $b_p \\leq 0.3$ since the base is so small that the filtering factor reaches almost immediately zero. The best orientation is found for $b_p = 0.7$.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.00\\columnwidth]{.\/figures\/op_b_p}\n\t\\caption{Overall algorithm results varying the $b_p$ operating parameter.}\n\t\\label{fig:opBP}\n\\end{figure}\n\n\\textbf{\\emph{Photograph edge filtering maximum segment length ($l_p$)}}: The choice of the maximum segment length for the photograph edge filtering is introduced to allow mountain edges which are thicker than one pixel not to penalize the other edges placed below it, but at the same time to not allow the noise edges that generate columns of points to be left unfiltered. The dependency of the result on $l_p$ is displayed in Figure \\ref{fig:opLP}: the optimal balance between the two trends is reached with $l_p = 2$.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.00\\columnwidth]{.\/figures\/op_l_p}\n\t\\caption{Overall algorithm results varying the $l_p$ operating parameter.}\n\t\\label{fig:opLP}\n\\end{figure}\n\n\\textbf{\\emph{Panorama edge filtering base ($b_r$)}}: The choice of the filtering base for the panorama edges defines how strongly the lower edge points will be reduced in strength, $b_r = 1.0$ means that no filtering will be performed, and $b_r = 0.0$ that only the first segment of each row of the photograph will be left. The dependency of the result on $b_p$ is displayed in Figure \\ref{fig:opBP}: the overall trend is similar to the filtering base of photograph edges ($b_p$) but with an important difference for high values: the result does not get worse with $b_r$ getting close to $1$. This can be explained by the fact that the high values ($\\geq 0.8$) of $b_p$ do not perform a significant filtering, so the noise edge points are not filtered and penalize the matching: the panorama edges instead do not have any noise, so do not have the same trend. The choice to set $b_r = 1.0$ is very clear, or in other words to not perform any filtering of panorama edge points.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.00\\columnwidth]{.\/figures\/op_b_r}\n\t\\caption{Overall algorithm results varying the $b_r$ operating parameter.}\n\t\\label{fig:opBR}\n\\end{figure}\n\n\\textbf{\\emph{Panorama edge filtering maximum segment length ($l_r$)}}: since $b_r$ is set to $1$ (the filtering of the panorama edges is not performed) the maximum segment length for the panorama edges filtering has no impact on the result, so its choice is not significant.\n\n\\textbf{\\emph{Photograph scaling interval ($\\pm k \\%$)}}: The choice of the scaling interval with respect to the originally estimated scaling process is described in the previous chapters. The results of its estimation were probably the most unexpected of all operating parameters. The dependency of the result on $b_p$ is displayed in Figure \\ref{fig:opK}: the best performance is reached with $k = 0$ (no scaling interval, only the estimated scale factor) and always worsens with increasing $k$. This trend can be interpreted as the proof of the fact that the Field of View and scaling factor estimation are so precise that increasing the scaling interval brings more penalties by setting an imprecise scale factor and increasing the probability of a noise edge matching the mountain panorama edge than advantages by trying different zoom factors and getting the best mountain edge overlap. The best orientation is therefore found with $k = 0$ (no scaling interval).\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.00\\columnwidth]{.\/figures\/op_k}\n\t\\caption{Overall algorithm results varying the $\\pm k \\%$ operating parameter.}\n\t\\label{fig:opK}\n\\end{figure}\n\\chapter{Conclusions and Future Work}\n\\label{Conclusions and Future Work}\n\\thispagestyle{empty}\n\n\\vspace{0.5cm}\n\n\\noindent\nIn this work an algorithm for the estimation of the orientation of a geo-tagged photograph containing mountains and for the identification of the visible mountain peaks is presented in detail. The algorithm has been implemented and evaluated with the result of estimating correctly the direction of the view of \\textbf{64.2\\%} of the various photographs.\n\nThis result can be considered excellent if used in passive crowdsourcing applications: given the enormous availability of photographs on the Web, the desired amount of correctly matched photographs can be reached just by increasing the number of processed photographs.\n\nIt is a very promising even if not excellent result speaking about applications connected with user experience: approximately one photograph out of three may not be matched, decreasing significantly the usability of the application.\n\n\\section{Future enhancements}\nWe expect to implement and test several techniques in the near future, aimed at enhancing the matching process and improving the results.\n\nOne of the main reasons that the matching fails is the massive presence of clouds on the photograph, and to prevent the edge filtering step from emphasizing the noise cloud edges a sky\/ground segmentation step can be inserted before the edge extraction phase: a sky\/ground segmentation algorithm (such as, for example, that proposed by Todorovic and Nechyba \\cite{Todorovic03sky\/groundmodeling}) that given a photograph splits it into the two regions of sky and terrain, then the sky part is erased from the photograph and the matching algorithm goes on without the cloud edge noise. In Figure \\ref{fig:skyTerrain} an example of this approach and the improvement it brings to the filtered edges map is shown.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/sky_terrain_enhancement}\n\t\\caption{An example of the use of sky\/ground segmentation in the edge extraction technique: noise cloud edges are removed and the mountain boundary is perfectly extracted.}\n\t\\label{fig:skyTerrain}\n\\end{figure}\n\nAnother frequent reason for incorrect matching is the imperfections between the boundaries on the photograph and on the render of mountains very close to the observer. Due to the small and unavoidable errors in the registered position, altitude estimation, and elevation model, the rendered panorama will always be imperfect with respect to reality, and this imperfection, for obvious reasons, get smaller as the distance from the observer to the object increases. The situation of having mountain boundaries both in the foreground and the background is very frequent: walking in a mountainous area the observer is usually surrounded by mountains placed close to him that are obscuring the distant mountains; the dimensions and the majesty of the mountains however bring usually the photographer to take pictures of distantly placed mountains with altitude higher that his point of view. For this reason as soon as the mountains in the foreground get placed in a way to create an \"aperture\", a photograph of the mountains visible in this sort of window will be probably shot. An example of this type of photograph can be that used for Figure \\ref{fig:skyTerrain} or that shown in Figure \\ref{fig:viewAperture}: it represents exactly the described situation, and it can be easily seen that the photograph and panorama edges of the mountain in the background are perfectly matchable, while the edges belonging to the closer mountains placed in the foreground are significantly different. In this case the algorithm manages to correctly alignment of the photograph anyway, but it is frequent in these cases of funnel-shaped mountain apertures that it fail, specially if the aperture is very small with respect to the total photograph dimension. Future studies are planned to include the development of a technique for the recognition of this type of situation and emphasizing the edges of the background mountains with respect to the foreground mountains.\n\n\\begin{figure}[h!]\n\t\\centering\n\t\\includegraphics[width=1.0\\columnwidth]{.\/figures\/view_aperture}\n\t\\caption{An example of a photograph (top) of background mountains seen in the aperture between the foreground mountains and the corresponding edge matching (bottom, red - panorama, blue - photograph).}\n\t\\label{fig:viewAperture}\n\\end{figure}\n\nEven when the estimated orientation is not correct, in most cases the correct alignment presents a significant peak in the VCC score distribution, so the best edge overlap estimation could be improved with a robust matching technique operating on the top-N positions extracted from the VCC matching score distribution. The likely method of implementing it is the neighborhood metrics of the edges as proposed by Baboud et al. \\cite{Baboud2011Alignment}; even if a simplified implementation of the robust matching technique proposed was implemented in this work and rejected as self-defeating, we believe that it can be refined to reach a higher rate of correctly matched photographs.\n\n\\section{Crowdsourcing involvement}\nAn important future direction of the research of this work is the integration with (traditional) crowdsourcing. Though being a feasible image processing task, photograph-to-panorama alignment is definitely better performed by humans that by algorithms. There is no need to move the photograph through the panorama or know the correct scale factor: a person can usually easily align the photograph to the panorama just by looking at it, comparing the fragments of the photograph he finds most particular, which can be elevated peaks, details that catch the eye and in general sense any terrain silhouette fragments that the person considers unusual. This allows one to easily identify the correct alignment by eye even in presence of significant errors between the photograph and the generated panorama.\n\nThis consideration leads us to the integration of the current work with crowdsourcing with several possible scenarios:\n\\begin{itemize}\n\\item\n\\emph{Contribution in learning and testing phase}:\nmanual photograph-to-panorama matching estimation can be very useful as a source of ground truth data both for the learning phase in case of using a machine intelligence algorithm in the edge matching and for the implementation testing phase, validating the results proposed by the algorithm. These approaches can be applied both to expert and non-expert crowdsourcing options: the non-expert tasks can consist of an activity feasible by anyone, such as searching for the correct alignment between the photograph and the corresponding panorama, the expert tasks instead aim at the activities that only mountaineers can perform, such as tagging the mountain peaks on a photograph based on his own knowledge and experience. Both methods allow the validation of the algorithm estimation result, but at two different (even if very related) levels: the first at photograph direction estimation, and the second at final mountain peak identification and tagging.\n\\item\n\\emph{Contribution as post processing validation}:\nhuman validation can be used not only in the development phase but also in the final implementation pipeline as the post processing validation. If the application is time-critical such as real-time augmented reality or website photographs mountain peak tagging, this approach cannot be applied, but in crawler applications such as environmental model creation, crowdsourcing can provide a significant improvement in the data quality by proposing manual photograph-to-panorama matching for the photographs that the algorithm has marked as photographs with low confidence score. Crowdsourcing therefore will be used as a complementary technique for photographs that cannot be aligned automatically.\n\\item\n\\emph{Contribution to data set expansion}:\nthe number of the available photographs in the data set is fundamental both for testing purposes and for the data to be processed by an environmental modeling system. Photographs can be collected from public sources such as social networks or photo sharing websites, filtering the geographical area and image content of the interest, or can be collected directly from people, uploading or signaling the photographs containing mountains in a certain requested area, or even photographs containing certain requested mountain peaks. This approach is important in the case of environmental modeling, when precisely collected photographs in the same area of ground truth data availability is fundamental.\n\\end{itemize}\n\n\\section{Possible areas of application}\nThere are several application areas the proposed algorithm can be used for, starting from the technique for snow level monitoring and prediction (the technique this work has been proposed and started for) and extending to possible applications in mobile and web fields that can be created thanks to this algorithm.\n\n\\subsection{Environmental Modeling}\nOne of the main future directions of the research will be the modeling of environmental processes by the analysis of mountain appearances extracted by the algorithm proposed in this work. The most important and most obvious measurements available from a mountain's appearance are the snow level and the snow water equivalent (SWE). The idea is to collect a series of photographs through time for each analyzed mountain, and based on the snow data ground truth, estimate the current measurements. The first step will therefore be an attempt to detect the correlation between the visual content of the mountain portion of a photograph with the physical measurements of the snow and SWE of that mountain. An interesting planned approach is to exploit also webcams in the region of interest since they present several advantages:\n\\begin{itemize}\n\\item\nweather and tourist webcams are very popular and frequent in mountain regions\n\\item\nthe time density of the measures can be as high as we want (it is only a matter of how frequently the image is acquired from the webcam, and in any case there is no reason to suppose the need of more than a couple of captures per day)\n\\item\nthe mountains captured on a webcam are always in the same position on the image, so even in cases of difficult edge matching it can be done or verified manually once to exploit all the instances of the photographs in future.\n\\end{itemize}\n\nFurthermore, supposing the mountains to be greatly distant from the observer (a reasonable and weak assumption in case of mountains photographs), by knowing the estimation of the observer's altitude and the altitude of the identified peak, we can estimate the altitude of each pixel of the mountain on the photograph by a simple proportion of the differences of altitudes and the pixel height of the mountain. This allows a comparison between the visual features of partial photographs corresponding to the same altitude even for photographs with significantly different shot position and photo camera properties.\n\n\\subsection{Augmented Reality}\nAugmented reality applications on mobile devices (applications that augment and supply the real-time camera view of the mobile device with computer-generated input) is a recent niche topic, and the promising use of the described algorithm regards augmented reality: an application can tag in real-time the mountains viewed by the user, highlight the peaks and terrain silhouettes, and augment the mountains with useful information such as altitude contours drawn on the image.\n\nSuch a kind of application would eliminate the problem of wrong geo-tag estimation (keeping the GPS of the mobile device on, the position is usually estimated with a tollerance of few meters) and so will reduce significantly the problem of wrong altitude estimation and elevation model imperfections with respect to reality. The reduced computation capacity may be compensated by the built-in compass, which gives a rough indication of the observer's direction of view, so the matching procedure can be done only on a reduced fragment of the rendered panorama. The bandwidth use will be small since the mountains are usually distant from the observer, so the rendered view will change very slowly while the observer is moving, so it will need to be updated rarely.\n\n\\subsection{Photo Sharing Platforms}\nTagging the mountain peaks on the geo-tagged photographs can lead also to a significant improvement to a cataloging and searching system of a social network or a photo sharing website, and to a better user experience by exploring the peak names and other information by directly viewing the photograph.\nAutomatic tagging of mountain peaks (tagging intended as the catalog assignment of peak names to the photograph and not the visual annotation on the photograph itself) can allow navigation through mountain photographs in these scenarios:\n\\begin{itemize}\n\\item\nUser searches for a mountain by specifying its name in the query: retrieves the photographs of that mountain, even if the author of the photograph did not specify the name in the title, description or other metadata.\n\\item\nUser views a photograph: other photographs of the same mountain (with the same or different facade) are suggested, even if the authors of both photographs did not specify the name in the title, description or other metadata.\n\\end{itemize}\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\nConsider a set of agents $\\mathcal{N}=\\{1,2,\\ldots,n\\}$ connected over a network. Each agent has a local smooth and strongly convex cost function $f_i:\\mathbb{R}^{p}\\rightarrow \\mathbb{R}$. The global objective is to locate $x\\in\\mathbb{R}^p$ that minimizes the average of all cost functions:\n\\begin{equation}\n\\min_{x\\in \\mathbb{R}^{p}}f(x)\\left(=\\frac{1}{n}\\sum_{i=1}^{n}f_i(x)\\right). \\label{opt Problem_def}\n\\end{equation}%\nScenarios in which problem (\\ref{opt Problem_def}) is considered include distributed machine learning \\cite{forrester2007multi,nedic2017fast,cohen2017projected,wai2018multi}, multi-agent target seeking \\cite{pu2016noise,chen2012diffusion}, and wireless networks \\cite{cohen2017distributed,mateos2012distributed,baingana2014proximal}, among many others.\n\nTo solve problem (\\ref{opt Problem_def}), we assume each agent $i$ queries a stochastic oracle ($\\mathcal{SO}$) to obtain noisy gradient samples of the form $g_i(x,\\xi_i)$ that satisfies the following condition:\n\\begin{assumption}\n\t\\label{asp: gradient samples}\n\tFor all $i\\in\\mathcal{N}$ and all $x\\in\\mathbb{R}^p$, \n\teach random vector $\\xi_i\\in\\mathbb{R}^m$ is independent, and\n\t\\begin{equation}\n\t\\begin{split}\n\t& \\mathbb{E}_{\\xi_i}[g_i(x,\\xi_i)\\mid x] = \\nabla f_i(x),\\\\\n\t& \\mathbb{E}_{\\xi_i}[\\|g_i(x,\\xi_i)-\\nabla f_i(x)\\|^2\\mid x] \\le \\sigma^2\\quad\\hbox{\\ for some $\\sigma>0$}.\n\t\\end{split}\n\t\\end{equation}\n\\end{assumption}\nThe above assumption of stochastic gradients holds true for many on-line distributed learning problems, where $f_i(x)=\\mathbb{E}_{\\xi_i}[F_i(x,\\xi_i)]$ denotes the expected loss function agent $i$ wishes to minimize, while independent samples $\\xi_i$ are gathered continuously over time.\nFor another example, in simulation-based optimization, the gradient estimation often incurs noise \nthat can be due to various sources, such as modeling and discretization errors, \nincomplete convergence, and finite sample size for Monte-Carlo methods~\\cite{kleijnen2008design}. \n\nDistributed algorithms dealing with problem (\\ref{opt Problem_def}) have been studied extensively in the literature \\cite{tsitsiklis1986distributed,nedic2009distributed,nedic2010constrained,jakovetic2014fast,kia2015distributed,shi2015extra,di2016next,qu2017harnessing,nedic2017achieving,pu2018push}.\nRecently, there has been considerable interest in distributed implementation of stochastic gradient algorithms (see \\cite{ram2010distributed,cavalcante2013distributed,towfic2014adaptive,lobel2011distributed,srivastava2011distributed,wang2015cooperative,lan2017communication,pu2017flocking,lian2017can,pu2018swarming}). The literature has shown that distributed algorithms may compete with, or even outperform, their centralized counterparts under certain conditions \\cite{pu2017flocking,lian2017can,pu2018swarming}. For instance, in our recent work~\\cite{pu2018swarming}, we proposed a swarming-based approach for distributed stochastic optimization which beats a centralized gradient method in real-time assuming that all $f_i$ are identical. \nHowever, to the best of our knowledge, there is no distributed stochastic gradient method addressing problem (\\ref{opt Problem_def}) that shows comparable performance with a centralized approach. In particular, under constant step size policies none of the existing algorithms achieve an error bound that is decreasing in the network size $n$. \n\n\nA distributed gradient tracking method was proposed in~\\cite{di2016next,nedic2017achieving,qu2017harnessing}, where \nthe agent-based auxiliary variables $y_i$ were introduced to track the average gradients of $f_i$ assuming accurate gradient information is available. It was shown that the method, with constant step size, generates iterates that converge linearly to the optimal solution.\nInspired by the approach, in this paper we consider a distributed stochastic gradient tracking method. By comparison, in our proposed algorithm $y_i$ are tracking the stochastic gradient averages of $f_i$.\nWe are able to show that the iterates generated by each agent reach, in expectation, a neighborhood of the optimal point exponentially fast under a constant step size.\nInterestingly, with a sufficiently small step size, the limiting error bounds on the distance \nbetween agent iterates and the optimal solution decrease in the network size $n$, which is comparable to the performance of a centralized stochastic gradient algorithm.\n\nOur work is also related to the extensive literature in stochastic approximation (SA) methods dating back to the seminal works~\\cite{robbins1951stochastic} and~\\cite{kiefer1952stochastic}. These works include the analysis of convergence (conditions for convergence, rates of convergence, suitable choice of step size) in the context of diverse noise models~\\cite{kushner2003stochastic}.\n\nThe paper is organized as follows. In Section~\\ref{sec: set}, we introduce the distributed stochastic gradient tracking method along with the main results. We perform analysis in Section~\\ref{sec:cohesion} \nand provide a numerical example in Section~\\ref{sec: simulation} to illustrate our theoretical findings. \nSection~\\ref{sec: conclusion} concludes the paper.\n\n\n\\subsection{Notation}\n\\label{subsec:pre}\nThroughout the paper, vectors default to columns if not otherwise specified.\nLet each agent $i$ hold a local copy $x_i\\in\\mathbb{R}^p$ of the decision variable and an auxiliary variable $y_i\\in\\mathbb{R}^p$. Their values at iteration\/time $k$ are denoted by $x_{i,k}$ and $y_{i,k}$, respectively. \nWe let\n\\begin{equation*}\n \\mathbf{x} := [x_1, x_2, \\ldots, x_n]^{\\intercal}\\in\\mathbb{R}^{n\\times p},\\ \\ \\mathbf{y} := [y_1, y_2, \\ldots, y_n]^{\\intercal}\\in\\mathbb{R}^{n\\times p},\n\\end{equation*}\n\\begin{equation}\n\\overline{x} := \\frac{1}{n}\\mathbf{1}^{\\intercal} \\mathbf{x}\\in\\mathbb{R}^{1\\times p},\\ \\ \\overline{y} := \\frac{1}{n}\\mathbf{1}^{\\intercal}\\mathbf{y}\\in\\mathbb{R}^{1\\times p},\n\\end{equation}\nwhere $\\mathbf{1}$ denotes the vector with all entries equal to 1.\nWe define an aggregate objective function of the local variables:\n\\begin{equation}\nF(\\mathbf{x}):=\\sum_{i=1}^nf_i(x_i),\n\\end{equation}\nand write\n\\begin{equation*}\n\\nabla F(\\mathbf{x}):=\\left[\\nabla f_1(x_1), \\nabla f_2(x_2), \\ldots, \\nabla f_n(x_n)\\right]^{\\intercal}\\in\\mathbb{R}^{n\\times p}.\n\\end{equation*}\nIn addition, let\n\\begin{equation}\n\\label{def: h}\nh(\\mathbf{x}):=\\frac{1}{n}\\mathbf{1}^{\\intercal}\\nabla F(\\mathbf{x})\\in\\mathbb{R}^{1\\times p},\n\\end{equation}\n\\begin{equation*}\n\\boldsymbol{\\xi} := [\\xi_1, \\xi_2, \\ldots, \\xi_n]^{\\intercal}\\in\\mathbb{R}^{n\\times p},\n\\end{equation*}\nand\n\\begin{equation}\nG(\\mathbf{x},\\boldsymbol{\\xi}):=[g_1(x_1,\\xi_1), g_2(x_2,\\xi_2), \\ldots, g_n(x_n,\\xi_n)]^{\\intercal}\\in\\mathbb{R}^{n\\times p}.\n\\end{equation}\n\nInner product of two vectors $a,b$ of the same dimension is written as $\\langle a,b\\rangle$. For two matrices $A,B\\in\\mathbb{R}^{n\\times p}$, we define\n\\begin{equation}\n\\langle A,B\\rangle :=\\sum_{i=1}^n\\langle A_i,B_i\\rangle,\n\\end{equation}\nwhere $A_i$ (respectively, $B_i$) represents the $i$-th row of $A$ (respectively, $B$). We use $\\|\\cdot\\|$ to denote the $2$-norm of vectors; for matrices, $\\|\\cdot\\|$ denotes the Frobenius norm.\n\nA graph is a pair $\\mathcal{G}=(\\mathcal{V},\\mathcal{E})$ where $\\mathcal{V}=\\{1,2,\\ldots,n\\}$ is the set of vertices (nodes) and $\\mathcal{E}\\subseteq \\mathcal{V}\\times \\mathcal{V}$ represents the set of edges connecting vertices. We assume agents communicate in an undirected graph, i.e., $(i,j)\\in\\mathcal{E}$ iff $(j,i)\\in\\mathcal{E}$.\nDenote by $W=[w_{ij}]\\in\\mathbb{R}^{n\\times n}$ the coupling matrix of agents. Agent $i$ and $j$ are connected iff $w_{ij}=w_{ji}>0$ ($w_{ij}=w_{ji}=0$ otherwise). Formally, we assume the following condition regarding the interaction among agents:\n\\begin{assumption}\n\t\\label{asp: network}\n\tThe graph $\\mathcal{G}$ corresponding to the network of agents is undirected and connected (there exists a path between any two agents). Nonnegative coupling matrix $W$ is doubly stochastic, \n\ti.e., $W\\mathbf{1}=\\mathbf{1}$ and $\\mathbf{1}^{\\intercal}W=\\mathbf{1}^{\\intercal}$.\n\tIn addition, $w_{ii}>0$ for some $i\\in\\mathcal{N}$.\n\\end{assumption}\nWe will frequently use the following result, which is a direct implication of Assumption \\ref{asp: network} (see \\cite{qu2017harnessing} Section II-B):\n\\begin{lemma}\n\t\\label{lem: spectral norm}\n\tLet Assumption \\ref{asp: network} hold, and let $\\rho_w$ denote the spectral norm of \n\tthe matrix $W-\\frac{1}{n}\\mathbf{1}\\mathbf{1}^{\\intercal}$. Then, $\\rho_w<1$ and \n\t\\begin{equation*}\n\t\\|W\\omega-\\mathbf{1}\\overline{\\omega}\\|\\le \\rho_w\\|\\omega-\\mathbf{1}\\overline{\\omega}\\|\n\t\\end{equation*}\n\tfor all $\\omega\\in\\mathbb{R}^{n\\times p}$, where $\\overline{\\omega}=\\frac{1}{n}\\mathbf{1}^{\\intercal}\\omega$.\n\\end{lemma}\n\n\\section{A Distributed Stochastic Gradient Tracking Method}\n\\label{sec: set}\nWe consider the following distributed stochastic gradient tracking method: at each step $k\\in\\mathbb{N}$, \nevery agent $i$ independently implements the following two steps:\n\\begin{equation}\n\\label{eq: x_i,k}\n\\begin{split}\nx_{i,k+1} = & \\sum_{j=1}^{n}w_{ij}(x_{j,k}-\\alpha y_{j,k}), \\\\\ny_{i,k+1} = & \\sum_{j=1}^{n}w_{ij}y_{j,k}+g_i(x_{i,k+1},\\xi_{i,k+1})-g_i(x_{i,k},\\xi_{i,k}),\n\\end{split}\n\\end{equation}\nwhere $\\alpha>0$ is a constant step size. The iterates are initiated with an arbitrary \n$x_{i,0}$ and $y_{i,0}= g_i(x_{i,0},\\xi_{i,0})$ for all~$i\\in{\\cal N}$.\nWe can also write (\\ref{eq: x_i,k}) in the following compact form:\n\\begin{equation}\n\\label{eq: x_k}\n\\begin{split}\n\\mathbf{x}_{k+1} = & W(\\mathbf{x}_k-\\alpha \\mathbf{y}_k), \\\\\n\\mathbf{y}_{k+1} = & W\\mathbf{y}_k+G(\\mathbf{x}_{k+1},\\boldsymbol{\\xi}_{k+1})-G(\\mathbf{x}_k,\\boldsymbol{\\xi}_k).\n\\end{split}\n\\end{equation}\nAlgorithm (\\ref{eq: x_i,k}) is closely related to the schemes considered in \\cite{di2016next,nedic2017achieving,qu2017harnessing}, where auxiliary variables $y_{i,k}$ \nwere introduced to track the average $\\frac{1}{n}\\sum_{i=1}^{n}\\nabla f_i(x_{i,k})$. This design ensures that the algorithm achieves linear convergence under constant step size choice.\nCorrespondingly, under our approach $y_{i,k}$ are (approximately) tracking $\\frac{1}{n}\\sum_{i=1}^{n}g_i(x_{i,k},\\xi_{i,k})$.\nTo see why this is the case, note that\n\\begin{equation}\n\\overline{y}_k = \\frac{1}{n}\\mathbf{1}^{\\intercal} \\mathbf{y}_k.\n\\end{equation}\nSince $y_{i,0}=g(x_{i,0},\\xi_{i,0}),\\forall i$. By induction,\n\\begin{equation}\n\\overline{y}_k=\\frac{1}{n}\\mathbf{1}^{\\intercal}G(\\mathbf{x}_k,\\boldsymbol{\\xi}_k),\\forall k.\n\\end{equation}\nWe will show that $\\mathbf{y}_k$ is close to $\\mathbf{1}\\overline{y}_k$ at each round. Hence $y_{i,k}$ are (approximately) tracking $\\frac{1}{n}\\sum_{i=1}^{n}g_i(x_{i,k},\\xi_{i,k})$.\n\nThroughout the paper, we make the following standing assumption on the objective functions $f_i$:\n\\begin{assumption}\n\t\\label{asp: strconvexity}\n\tEach $f_i:\\mathbb{R}^p\\rightarrow \\mathbb{R}$ is $\\mu$-strongly convex with $L$-Lipschitz continuous gradients, i.e., for any $x,x'\\in\\mathbb{R}^p$,\n\t\\begin{equation}\n\t\\begin{split}\n\t& \\langle \\nabla f_i(x)-\\nabla f_i(x'),x-x'\\rangle\\ge \\mu\\|x-x'\\|^2,\\\\\n\t& \\|\\nabla f_i(x)-\\nabla f_i(x')\\|\\le L \\|x-x'\\|.\n\t\\end{split}\n\t\\end{equation}\n\\end{assumption}\nUnder Assumption \\ref{asp: strconvexity}, problem (\\ref{opt Problem_def}) has a unique solution denoted by $x^*\\in\\mathbb{R}^{1\\times p}$.\n\n\\subsection{Main Results}\n\\label{subsec: main}\nMain convergence properties of the distributed gradient tracking method (\\ref{eq: x_i,k}) are covered in the following theorem.\n\\begin{theorem}\n\t\\label{Theorem}\n\tSuppose Assumptions \\ref{asp: gradient samples}-\\ref{asp: strconvexity} hold and $\\alpha$ satisfies\n\t\\begin{equation}\n\t\t\\label{alpha_ultimate_bound}\n\t\t\\alpha\\le \\min\\left\\{\\frac{(1-\\rho_w^2)}{12\\rho_w^2 L},\\frac{(1-\\rho_w^2)^2}{2\\sqrt{\\Gamma}L\\max\\{6\\rho_w\\|W-I\\|,1-\\rho_w^2\\}}, \\frac{(1-\\rho_w^2)}{3\\rho_w^{2\/3}L}\\left[\\frac{\\mu^2}{L^2}\\frac{(\\Gamma-1)}{\\Gamma(\\Gamma+1)}\\right]^{1\/3}\\right\\}\n\t\t\\end{equation}\n\tfor some $\\Gamma>1$. Then both $\\sup_{l\\ge k}\\mathbb{E}[\\|\\overline{x}_l-x^*\\|^2]$ and $\\sup_{l\\ge k}\\mathbb{E}[\\|\\mathbf{x}_{l+1}-\\mathbf{1}\\overline{x}_{l+1}\\|^2]$ converge at the linear rate $\\mathcal{O}(\\rho_A^k)$, where $\\rho_A<1$ is the spectral radius of\n\t\\begin{eqnarray*}\n\t\t\tA=\\begin{bmatrix}\n\t\t\t\t1-\\alpha\\mu & \\frac{\\alpha L^2}{\\mu n}(1+\\alpha\\mu) & 0\\\\\n\t\t\t\t0 & \\frac{1}{2}(1+\\rho_w^2) & \\alpha^2\\frac{(1+\\rho_w^2)\\rho_w^2}{(1-\\rho_w^2)}\\\\\n\t\t\t\t2\\alpha nL^3 & \\left(\\frac{1}{\\beta}+2\\right)\\|W-I\\|^2 L^2+3\\alpha L^3 & \\frac{1}{2}(1+\\rho_w^2)\n\t\t\t\\end{bmatrix},\n\t\\end{eqnarray*}\nin which $\\beta=\\frac{1-\\rho_w^2}{2\\rho_w^2}-4\\alpha L-2\\alpha^2 L^2$.\nFurthermore,\n\\begin{equation}\n\t\\label{error_bound_ultimate}\n\t\\limsup_{k\\rightarrow\\infty}\\mathbb{E}[\\|\\overline{x}_k-x^*\\|^2]\\le \\frac{(\\Gamma+1)}{\\Gamma}\\frac{\\alpha\\sigma^2}{\\mu n}\n\t+\\left(\\frac{\\Gamma+1}{\\Gamma-1}\\right)\\frac{4\\alpha^2 L^2(1+\\alpha\\mu)(1+\\rho_w^2)\\rho_w^2}{\\mu^2 n(1-\\rho_w^2)^3}M_{\\sigma},\n\t\\end{equation}\nand\n\\begin{equation}\n\t\\label{consensus_error_bound_ultimate}\n\t\\limsup_{k\\rightarrow\\infty}\\mathbb{E}[\\|\\mathbf{x}_k-\\mathbf{1}\\overline{x}_k\\|^2]\n\t\\le \\left(\\frac{\\Gamma+1}{\\Gamma-1}\\right)\\frac{4\\alpha^2 (1+\\rho_w^2)\\rho_w^2(2\\alpha^2L^3\\sigma^2+\\mu M_{\\sigma})}{\\mu(1-\\rho_w^2)^3},\n\t\\end{equation}\nwhere \n\\begin{equation}\n\t\\label{M_sigma}\n\tM_{\\sigma}:=\\left[3\\alpha^2L^2+2(\\alpha L+1)(n+1)\\right]\\sigma^2.\n\t\\end{equation}\n\\end{theorem}\n\\begin{remark}\n\tThe first term on the right hand side of (\\ref{error_bound_ultimate}) can be interpreted as the error caused by stochastic gradients only, since it does not depend on the network topology. The second term as well as the bound in (\\ref{consensus_error_bound_ultimate}) are network dependent and increase with $\\rho_w$ (larger $\\rho_w$ indicates worse network connectivity).\n\t\n\tIn light of (\\ref{error_bound_ultimate}) and (\\ref{consensus_error_bound_ultimate}),\n\t\\begin{equation*}\n\t\\limsup_{k\\rightarrow\\infty}\\mathbb{E}[\\|\\overline{x}_k-x^*\\|^2]=\\alpha\\mathcal{O}\\left(\\frac{\\sigma^2}{\\mu n}\\right)+\\alpha^2\\mathcal{O}\\left(\\frac{ L^2\\sigma^2}{\\mu^2}\\right),\n\t\\end{equation*}\n\tand\n\t\\begin{equation*}\n\t\\limsup_{k\\rightarrow\\infty}\\frac{1}{n}\\mathbb{E}[\\|\\mathbf{x}_k-\\mathbf{1}\\overline{x}_k\\|^2]=\\alpha^4\\mathcal{O}\\left(\\frac{ L^3\\sigma^2}{\\mu n}\\right)+\\alpha^2\\mathcal{O}\\left(\\sigma^2\\right).\n\t\\end{equation*}\n\tLet $(1\/n)\\mathbb{E}[\\|\\mathbf{x}_k-\\mathbf{1}x^*\\|^2]$ measure the average quality of solutions obtained by all the agents. We have\n\t\t\\begin{equation*}\n\t\\limsup_{k\\rightarrow\\infty}\\frac{1}{n}\\mathbb{E}[\\|\\mathbf{x}_k-\\mathbf{1}x^*\\|^2]=\\alpha\\mathcal{O}\\left(\\frac{\\sigma^2}{\\mu n}\\right)+\\alpha^2\\mathcal{O}\\left(\\frac{ L^2\\sigma^2}{\\mu^2}\\right),\n\t\\end{equation*}\n\twhich is decreasing in the network size $n$ when $\\alpha$ is sufficiently small\\footnote{Although $\\rho_w$ is also related to the network size $n$, it only appears in the terms with high orders of $\\alpha$.}.\n\t\n\tUnder a centralized algorithm in the form of\n\t\\begin{equation}\n\t\\label{eq: centralized}\n\tx_{k+1}=x_k-\\alpha \\frac{1}{n}\\sum_{i=1}^n g_i(x_k,\\xi_{i,k}),\\ \\ k\\in\\mathbb{N},\n\t\\end{equation}\n\twe would obtain\n\t\\begin{equation*}\n\\limsup_{k\\rightarrow\\infty}\\mathbb{E}[\\|x_k-x^*\\|^2]=\\alpha\\mathcal{O}\\left(\\frac{\\sigma^2}{\\mu n}\\right).\n\t\\end{equation*}\n\tIt can be seen that the distributed stochastic gradient tracking method (\\ref{eq: x_i,k}) is comparable with the centralized algorithm (\\ref{eq: centralized}) in their ultimate error bounds (up to constant factors) with sufficiently small step sizes.\n\\end{remark}\n\n\\begin{corollary}\n\t\\label{cor: speed}\n\tUnder the conditions in Theorem {\\ref{Theorem}}. Suppose in addition\\footnote{This condition is weaker than (\\ref{alpha_ultimate_bound}) in most cases.} \n\t\\begin{equation}\n\t\t\\label{alpha condition corollary}\n\t\\alpha\\le \\frac{(\\Gamma+1)}{\\Gamma}\\frac{(1-\\rho_w^2)}{8\\mu}.\n\t\\end{equation}\n Then\n\t\\begin{equation*}\n\t\\rho_A\\le 1-\\left(\\frac{\\Gamma-1}{\\Gamma+1}\\right)\\alpha\\mu.\n\t\\end{equation*}\n\\end{corollary}\n\\begin{remark}\n\tCorollary \\ref{cor: speed} implies that, for sufficiently small step sizes, \n\tthe distributed gradient tracking method has a comparable convergence speed to that of a centralized scheme (in which case the linear rate is $\\mathcal{O}(1-2\\alpha\\mu)^k$).\n\\end{remark}\n\n\\section{Analysis}\n\\label{sec:cohesion}\n\nIn this section, we prove Theorem \\ref{Theorem} by studying the evolution of $\\mathbb{E}[\\|\\overline{x}_k-x^*\\|^2]$, $\\mathbb{E}[\\|\\mathbf{x}_k-\\mathbf{1}\\overline{x}_k\\|^2]$ and $\\mathbb{E}[\\|\\mathbf{y}_k-\\mathbf{1}\\overline{y}_k\\|^2]$. Our strategy is to bound the three expressions in terms of linear combinations of their past values, in which way we establish a linear system of inequalities. This approach is different from those employed in \\cite{qu2017harnessing,nedic2017achieving}, where the analyses pertain to the examination of $\\|\\overline{x}_k-x^*\\|$, $\\|\\mathbf{x}_k-\\mathbf{1}\\overline{x}_k\\|$ and $\\|\\mathbf{y}_k-\\mathbf{1}\\overline{y}_k\\|$. Such distinction is due to the stochastic gradients $g_i(x_{i,k},\\xi_{i,k})$ whose variances play a crucial role in deriving the main inequalities.\n\nWe first introduce some lemmas that will be used later in the analysis. Denote by $\\mathcal{H}_k$ the history sequence $\\{\\mathbf{x}_0,\\boldsymbol{\\xi}_0,\\mathbf{y}_0,\\ldots,\\mathbf{x}_{k-1},\\boldsymbol{\\xi}_{k-1},\\mathbf{y}_{k-1},\\mathbf{x}_k\\}$, and define $\\mathbb{E}[\\cdot \\mid\\mathcal{H}_k]$ as the conditional expectation given $\\mathcal{H}_k$.\n\\begin{lemma}\n\t\\label{lem: oy_k-h_k}\n\tUnder Assumption \\ref{asp: gradient samples},\n\t\\begin{align}\n\t\\mathbb{E}\\left[\\|\\overline{y}_k-h(\\mathbf{x}_k)\\|^2\\mid\\mathcal{H}_k\\right] \\le \\frac{\\sigma^2}{n}.\n\t\\end{align}\n\\end{lemma}\n\\begin{proof}\n\tBy the definitions of $\\overline{y}_k$ and $h(\\mathbf{x}_k)$,\n\t\\begin{equation*}\n\t\t\\mathbb{E}\\left[\\|\\overline{y}_k-h(\\mathbf{x}_k)\\|^2\\mid\\mathcal{H}_k\\right]\\\\\n\t\t=\t\\frac{1}{n^2}\\sum_{i=1}^n\\mathbb{E}\\left[\\|g_i(x_{i,k},\\xi_{i,k})-\\nabla f_i(x_{i,k})\\|^2\\vert\\mathcal{H}_k\\right]\\le \\frac{\\sigma^2}{n}.\n\t\t\\end{equation*}\n\\end{proof}\n\\begin{lemma}\n\t\\label{lem: strong_convexity}\n\tUnder Assumption \\ref{asp: strconvexity},\n\t\\begin{align}\n\t\\| \\nabla F(\\overline{x}_k)-h(\\mathbf{x}_k)\\| \\le \\frac{L}{\\sqrt{n}}\\|\\mathbf{x}_k-\\mathbf{1}\\overline{x}_k\\|.\n\t\\end{align}\n\tSuppose in addition $\\alpha<2\/(\\mu+L)$. Then\n\t\\begin{equation*}\n\t\\|x-\\alpha\\nabla f(x)-x^*\\|\\le (1-\\alpha \\mu)\\|x-x^*\\|,\\, \\forall x\\in\\mathbb{R}^p.\n\t\\end{equation*}\n\\end{lemma}\n\\begin{proof}\n\tSee \\cite{qu2017harnessing} Lemma 10 for reference.\n\\end{proof}\n\nIn the following lemma, we establish bounds on $\\|\\mathbf{x}_{k+1}-\\mathbf{1}\\overline{x}_{k+1}\\|^2$ and on the conditional expectations of $\\|\\overline{x}_{k+1}-x^*\\|^2$ and $\\|\\mathbf{y}_{k+1}-\\mathbf{1}\\overline{y}_{k+1}\\|^2$, respectively.\n\\begin{lemma}\n\t\\label{lem: Main_Inequalities}\n\tSuppose Assumptions \\ref{asp: gradient samples}-\\ref{asp: strconvexity} hold and $\\alpha<2\/(\\mu+L)$. We have the following inequalities:\n\\begin{equation}\n\t\t\\label{First_Main_Inequality}\n\t\\mathbb{E}[\\|\\overline{x}_{k+1}-x^*\\|^2\\mid \\mathcal{H}_k]\n\t\\le \\left(1-\\alpha\\mu\\right)\\|\\overline{x}_k-x^*\\|^2\\\\\n\t+\\frac{\\alpha L^2}{\\mu n}\\left(1+\\alpha\\mu\\right)\\| \\mathbf{x}_k-\\mathbf{1}\\overline{x}_k\\|^2\n\t+\\frac{\\alpha^2\\sigma^2}{n},\n\\end{equation}\n\t\\begin{equation}\n\\label{Second_Main_Inequality}\n\\|\\mathbf{x}_{k+1}-\\mathbf{1}\\overline{x}_{k+1}\\|^2\n\\le \\frac{(1+\\rho_w^2)}{2}\\|\\mathbf{x}_k-\\mathbf{1}\\overline{x}_k\\|^2\\\\+\\alpha^2\\frac{(1+\\rho_w^2)\\rho_w^2}{(1-\\rho_w^2)}\\|\\mathbf{y}_{k}-\\mathbf{1}\\overline{y}_k\\|^2,\n\\end{equation}\nand for any $\\beta>0$,\n\\begin{multline}\n\\label{Third_Main_Inquality}\n\\mathbb{E}[\\|\\mathbf{y}_{k+1}-\\mathbf{1}\\overline{y}_{k+1}\\|^2\\mid \\mathcal{H}_k]\n\\le \\left(1+4\\alpha L+2\\alpha^2 L^2+\\beta\\right)\\rho_w^2\\mathbb{E}[\\|\\mathbf{y}_{k}-\\mathbf{1}\\overline{y}_{k}\\|^2\\mid \\mathcal{H}_k]\\\\\n+\\left(\\frac{1}{\\beta}\\|W-I\\|^2 L^2+2\\|W-I\\|^2 L^2+3\\alpha L^3\\right)\\|\\mathbf{x}_k-\\mathbf{1}\\overline{x}_k\\|^2+2\\alpha nL^3\\|\\overline{x}_k-x^*\\|^2+M_{\\sigma}.\n\\end{multline}\n\\end{lemma}\n\\begin{proof}\n\tSee Appendix \\ref{proof lem: Main_Inequalities}.\n\t\\end{proof}\n\n\\subsection{Proof of Theorem \\ref{Theorem}}\n\tTaking full expectation on both sides of (\\ref{First_Main_Inequality}), (\\ref{Second_Main_Inequality}) and (\\ref{Third_Main_Inquality}), we obtain the following linear system of inequalities\n\\begin{eqnarray}\n\\label{linear_system}\n\\begin{bmatrix}\n\\mathbb{E}[\\|\\overline{x}_{k+1}-x^*\\|^2]\\\\\n\\mathbb{E}[\\|\\mathbf{x}_{k+1}-\\mathbf{1}\\overline{x}_{k+1}\\|^2]\\\\\n\\mathbb{E}[\\|\\mathbf{y}_{k+1}-\\mathbf{1}\\overline{y}_{k+1}\\|^2]\n\\end{bmatrix}\n\\le \nA\\begin{bmatrix}\n\\mathbb{E}[\\|\\overline{x}_{k}-x^*\\|^2]\\\\\n\\mathbb{E}[\\|\\mathbf{x}_{k}-\\mathbf{1}\\overline{x}_{k}\\|^2]\\\\\n\\mathbb{E}[\\|\\mathbf{y}_{k}-\\mathbf{1}\\overline{y}_{k}\\|^2]\n\\end{bmatrix}+\\begin{bmatrix}\n\\frac{\\alpha^2\\sigma^2}{n}\\\\\n0\\\\\nM_{\\sigma}\n\\end{bmatrix},\n\\end{eqnarray}\nwhere the inequality is to be taken component-wise, and the entries of the matrix\n$A=[a_{ij}]$ are given by\n\t\\begin{eqnarray*}\n\t\t& \\begin{bmatrix}\n\t\t\ta_{11}\\\\\n\t\t\ta_{21}\\\\\n\t\t\ta_{31}\n\t\t\\end{bmatrix} = \n\t\t\\begin{bmatrix}\n\t\t\t1-\\alpha\\mu\\\\\n\t\t\t0\\\\\n\t\t\t2\\alpha nL^3 \n\t\t\\end{bmatrix},\n\t\t\\begin{bmatrix}\n\t\t\ta_{12}\\\\\n\t\t\ta_{22}\\\\\n\t\t\ta_{32}\n\t\t\\end{bmatrix} = \n\t\t\\begin{bmatrix}\n\t\t\t\\frac{\\alpha L^2}{\\mu n}(1+\\alpha\\mu)\\\\\n\t\t\t\\frac{1}{2}(1+\\rho_w^2)\\\\\n\t\t\t\\left(\\frac{1}{\\beta}+2\\right)\\|W-I\\|^2 L^2+3\\alpha L^3\n\t\t\\end{bmatrix},\\\\\n\t\t& \\begin{bmatrix}\n\t\t\ta_{13}\\\\\n\t\t\ta_{23}\\\\\n\t\t\ta_{33}\n\t\t\\end{bmatrix} = \n\t\t\\begin{bmatrix}\n\t\t\t0 \\\\\n\t\t\t\\alpha^2\\frac{(1+\\rho_w^2)\\rho_w^2}{(1-\\rho_w^2)}\\\\\n\t\t\t\\left(1+4\\alpha L+2\\alpha^2 L^2+\\beta\\right)\\rho_w^2\n\t\t\\end{bmatrix},\n\\end{eqnarray*}\nand $M_{\\sigma}$ is given in (\\ref{M_sigma}).\nHence $\\sup_{l\\ge k}\\mathbb{E}[\\|\\overline{x}_l-x^*\\|^2]$, $\\sup_{l\\ge k}\\mathbb{E}[\\|\\mathbf{x}_l-\\mathbf{1}\\overline{x}_l\\|^2]$ and $\\sup_{l\\ge k}\\mathbb{E}[\\|\\mathbf{y}_l-\\mathbf{1}\\overline{y}_l\\|^2]$ all converge to a neighborhood of $0$ at the linear rate $\\mathcal{O}(\\rho_A^k)$ if the spectral radius of $A$ satisfies $\\rho_A<1$. \nThe next lemma provides conditions for relation $\\rho_A<1$ to hold.\n\\begin{lemma}\n\t\\label{lem: rho_M}\n\tLet $M=[m_{ij}]\\in\\mathbb{R}^{3\\times 3}$ be a nonnegative, irreducible matrix with \n\t$m_{ii}<\\lambda^*$ for some~{$\\lambda^*>0$ and all $i=1,2,3.$} \n\tA necessary and sufficient condition for $\\rho_M<\\lambda^*$ is $\\text{det}(\\lambda^* I-M)>0$.\n\\end{lemma}\n\\begin{proof}\n\tSee Appendix \\ref{subsec: proof lemma rho_M}.\n\t\\end{proof}\n\tLet $\\alpha$ and $\\beta$ be such that the following relations hold\\footnote{Matrix $A$ in Theorem~\\ref{Theorem} \n\tcorresponds to such a choice of $\\alpha$ and $\\beta$.}.\n\\begin{equation}\n\t\\label{beta}\n\ta_{33}=\\left(1+4\\alpha L+2\\alpha^2 L^2+\\beta\\right)\\rho_w^2=\\frac{1+\\rho_w^2}{2}<1,\n\t\\end{equation}\n\t\\begin{equation}\n\t\\label{alpha,beta}\n\ta_{23}a_{32}=\\alpha^2\\frac{(1+\\rho_w^2)\\rho_w^2}{(1-\\rho_w^2)}\\left[\\left(\\frac{1}{\\beta}+2\\right)\\|W-I\\|^2 L^2+3\\alpha L^3\\right]\n\t\\le\\frac{1}{\\Gamma}(1-a_{22})(1-a_{33})\n\t\\end{equation}\nfor some $\\Gamma>1$, and\n\\begin{equation}\n\t\\label{alpha last condition}\n\ta_{12}a_{23}a_{31}=\\frac{2\\alpha^4 L^5(1+\\alpha\\mu)}{\\mu}\\frac{(1+\\rho_w^2)}{(1-\\rho_w^2)}\\rho_w^2\n\t\\le \\frac{1}{\\Gamma+1}(1-a_{11})[(1-a_{22})(1-a_{33})-a_{23}a_{32}].\n\t\\end{equation}\nThen,\n\\begin{multline*}\n\t\\text{det}(I-A)=(1-a_{11})(1-a_{22})(1-a_{33})-(1-a_{11})a_{23}a_{32}-a_{12}a_{23}a_{31}\\\\\n\t\\ge \\frac{\\Gamma}{(\\Gamma+1)}(1-a_{11})[(1-a_{22})(1-a_{33})-a_{23}a_{32}]\n\t\\ge \\left(\\frac{\\Gamma-1}{\\Gamma+1}\\right)(1-a_{11})(1-a_{22})(1-a_{33})>0,\n\t\\end{multline*}\ngiven that $a_{11},a_{22},a_{33}<1$. In light of Lemma \\ref{lem: rho_M}, $\\rho_A<1$.\nIn addition, denoting $B:=[\\frac{\\alpha^2\\sigma^2}{n}, 0, M_{\\sigma}]^{\\intercal}$, we get \n\t\\begin{align}\n\t\t\\label{error_bound_preliminary}\n\t\\lim\\sup_{k\\rightarrow\\infty}\\mathbb{E}[\\|\\overline{x}_k-x^*\\|^2]\\le & [(I-A)^{-1}B]_1 \\notag\\\\\n\t= & \\frac{1}{\\text{det}(I-A)}\\left\\{\\left[(1-a_{22})(1-a_{33})-a_{23}a_{32}\\right]\\frac{\\alpha^2\\sigma^2}{n}+a_{12}a_{23}M_{\\sigma}\\right\\} \\notag\\\\\n\t\\le & \\frac{(\\Gamma+1)}{\\Gamma}\\frac{\\alpha\\sigma^2}{\\mu n}+\\left(\\frac{\\Gamma+1}{\\Gamma-1}\\right)\\frac{ a_{12}a_{23}M_{\\sigma}}{(1-a_{11})(1-a_{22})(1-a_{33})}\\notag\\\\\n\t= & \\frac{(\\Gamma+1)}{\\Gamma}\\frac{\\alpha\\sigma^2}{\\mu n}+ \\left(\\frac{\\Gamma+1}{\\Gamma-1}\\right)\\frac{\\alpha^3 L^2(1+\\alpha\\mu)(1+\\rho_w^2)\\rho_w^2M_{\\sigma}}{\\mu n(1-\\rho_w^2)(1-a_{11})(1-a_{22})(1-a_{33})}\\notag\\\\\n\t= & \\frac{(\\Gamma+1)}{\\Gamma}\\frac{\\alpha\\sigma^2}{\\mu n}+\\left(\\frac{\\Gamma+1}{\\Gamma-1}\\right)\\frac{4\\alpha^2 L^2(1+\\alpha\\mu)(1+\\rho_w^2)\\rho_w^2}{\\mu^2 n(1-\\rho_w^2)^3}M_{\\sigma},\n\t\\end{align}\nand\n\t\\begin{multline*}\n\t\\lim\\sup_{k\\rightarrow\\infty}\\mathbb{E}[\\|\\mathbf{x}_k-\\mathbf{1}\\overline{x}_k\\|^2] \\le [(I-A)^{-1}B]_2\n\t=\\frac{1}{\\text{det}(I-A)}\\left[a_{23}a_{31}\\frac{\\alpha^2\\sigma^2}{n}+a_{23}(1-a_{11})M_{\\sigma}\\right]\\\\\n\t\\le \\left(\\frac{\\Gamma+1}{\\Gamma-1}\\right)\\frac{a_{23}}{(1-a_{11})(1-a_{22})(1-a_{33})}\\left(2\\alpha nL^3\\frac{\\alpha^2 \\sigma^2}{n}+\\alpha\\mu M_{\\sigma}\\right)\\\\\n\t= \\frac{4(\\Gamma+1)\\alpha^2 (1+\\rho_w^2)\\rho_w^2(2\\alpha^2L^3\\sigma^2+\\mu M_{\\sigma})}{(\\Gamma-1)\\mu(1-\\rho_w^2)^3}.\n\t\\end{multline*}\n\n\tWe now show that (\\ref{beta}), (\\ref{alpha,beta}) and (\\ref{alpha last condition}) are satisfied under condition (\\ref{alpha_ultimate_bound}). By (\\ref{beta}) and (\\ref{alpha,beta}), it follows that\n\t\\begin{equation*}\n\t\t\\beta=\\frac{1-\\rho_w^2}{2\\rho_w^2}-4\\alpha L-2\\alpha^2 L^2>0,\n\t\t\\end{equation*}\n\tand\n\t\t\\begin{equation*}\n\t\t\\alpha^2\\frac{(1+\\rho_w^2)\\rho_w^2}{(1-\\rho_w^2)}\\left[\\left(\\frac{1}{\\beta}+2\\right)\\|W-I\\|^2 L^2+3\\alpha L^3\\right]\\le\\frac{(1-\\rho_w^2)^2}{4\\Gamma}.\n\t\t\\end{equation*}\n\t\tSince by (\\ref{alpha_ultimate_bound}) we have\n\t\t\\begin{equation*}\n\t\t\t\\alpha\\le \\frac{1-\\rho_w^2}{12\\rho_w^2 L},\\qquad\n\t\t\n\t\n\t\t\\beta\\ge \\frac{1-\\rho_w^2}{8\\rho_w^2}>0,\n\t\t\\end{equation*}\n\t\twe only need to show that\n\t\t\\begin{equation*}\n\t\t\\alpha^2\\left[\\frac{(2+6\\rho_w^2)}{(1-\\rho_w^2)}\\|W-I\\|^2 L^2+\\frac{(1-\\rho_w^2)}{4\\rho_w^2}L^2\\right]\\le\\frac{(1-\\rho_w^2)^3}{4\\Gamma(1+\\rho_w^2)\\rho_w^2}.\n\t\t\\end{equation*}\n\t\tThe preceding inequality is equivalent to\n\t\t\\begin{equation*}\n\t\t\\alpha\\le \\frac{(1-\\rho_w^2)^2}{L\\sqrt{\\Gamma(1+\\rho_w^2)}\\sqrt{4\\rho_w^2(2+6\\rho_w^2)\\|W-I\\|^2+(1-\\rho_w^2)^2}},\n\t\t\\end{equation*}\n\timplying that it is sufficient to have\n\t\\begin{equation*}\n\t\\alpha\\le\\frac{(1-\\rho_w^2)^2}{2\\sqrt{\\Gamma}L\\max(6\\rho_w\\|W-I\\|,1-\\rho_w^2)}.\n\t\\end{equation*}\n\t\tTo see that relation (\\ref{alpha last condition}) holds, consider a stronger condition\n\t\t\\begin{equation*}\n\t\t\\frac{2\\alpha^4 L^5(1+\\alpha\\mu)}{\\mu}\\frac{(1+\\rho_w^2)}{(1-\\rho_w^2)}\\rho_w^2\n\t\t\\le \\frac{(\\Gamma-1)}{\\Gamma(\\Gamma+1)}(1-a_{11})(1-a_{22})(1-a_{33}),\n\t\t\\end{equation*}\n\t\tor equivalently,\n\t\t\\begin{equation*}\n\t\t\\frac{2\\alpha^3 L^5(1+\\alpha\\mu)}{\\mu^2}\\frac{(1+\\rho_w^2)}{(1-\\rho_w^2)}\\rho_w^2 \\le \\frac{(\\Gamma-1)}{4\\Gamma(\\Gamma+1)}(1-\\rho_w^2)^2.\n\t\t\\end{equation*}\n\t\tIt suffices that\n\t\t\\begin{equation}\n\t\t\\alpha\\le \\frac{(1-\\rho_w^2)}{3\\rho_w^{2\/3}L}\\left[\\frac{\\mu^2}{L^2}\\frac{(\\Gamma-1)}{\\Gamma(\\Gamma+1)}\\right]^{1\/3}.\n\t\t\\end{equation}\n\n\\subsection{Proof of Corollary \\ref{cor: speed}}\nWe derive an upper bound of $\\rho_A$ under condition (\\ref{alpha_ultimate_bound}) and (\\ref{alpha condition corollary}). Note that the characteristic function of $A$ is given by\n\\begin{equation*}\n\\text{det}(\\lambda I-A)=(\\lambda-a_{11})(\\lambda-a_{22})(\\lambda-a_{33})\n-(\\lambda-a_{11})a_{23}a_{32}-a_{12}a_{23}a_{31}.\n\\end{equation*}\nSince $\\text{det}(I-A)> 0$ and $\\text{det}(\\max\\{a_{11},a_{22},a_{33}\\} I-A)<0$, $\\rho_A\\in(\\max\\{a_{11},a_{22},a_{33}\\},1)$. By (\\ref{alpha,beta}) and (\\ref{alpha last condition}),\n\\begin{multline*}\n\\text{det}(\\lambda I-A)\n\\ge (\\lambda-a_{11})(\\lambda-a_{22})(\\lambda-a_{33})-(\\lambda-a_{11})a_{23}a_{32}-\\frac{1}{\\Gamma+1}(1-a_{11})[(1-a_{22})(1-a_{33})-a_{23}a_{32}]\\\\\n\\ge (\\lambda-a_{11})(\\lambda-a_{22})(\\lambda-a_{33})-\\frac{1}{\\Gamma}(\\lambda-a_{11})(1-a_{22})(1-a_{33})\n-\\frac{(\\Gamma-1)}{\\Gamma(\\Gamma+1)}(1-a_{11})(1-a_{22})(1-a_{33}).\n\\end{multline*}\nSuppose $\\lambda=1-\\epsilon$ for some $\\epsilon\\in(0,\\alpha\\mu)$, satisfying\n\\begin{equation*}\n\\text{det}(\\lambda I-A)\n\\ge \\frac{1}{4}(\\alpha\\mu-\\epsilon)\\left(1-\\rho_w^2-2\\epsilon\\right)^2\n-\\frac{1}{4\\Gamma}(\\alpha\\mu-\\epsilon)(1-\\rho_w^2)^2-\\frac{(\\Gamma-1)\\alpha\\mu}{4\\Gamma(\\Gamma+1)}(1-\\rho_w^2)^2\\ge 0.\n\\end{equation*}\nUnder (\\ref{alpha condition corollary}), it suffices that\n\\begin{equation*}\n\\epsilon\\le \\left(\\frac{\\Gamma-1}{\\Gamma+1}\\right)\\alpha\\mu.\n\\end{equation*}\nDenote\n\\begin{equation*}\n\\tilde{\\lambda}=1-\\left(\\frac{\\Gamma-1}{\\Gamma+1}\\right)\\alpha\\mu.\n\\end{equation*}\nThen $\\text{det}(\\tilde{\\lambda} I-A)\\ge 0$ so that $\\rho_A\\le \\tilde{\\lambda}$.\n\n\\section{Numerical Example}\n\\label{sec: simulation}\n\nIn this section, we provide a numerical example to illustrate our theoretic findings. \nConsider the \\emph{on-line} Ridge regression problem, i.e.,\n\\begin{equation}\n\\label{Ridge Regression}\nf(x):=\\frac{1}{n}\\sum_{i=1}^n\\mathbb{E}_{u_i,v_i}\\left[\\left(u_i^{\\intercal} x-v_i\\right)^2+\\rho\\|x\\|^2\\right].\n\\end{equation}\nwhere $\\rho>0$ ia a penalty parameter.\nFor each agent $i$, samples in the form of $(u_i,v_i)$ are gathered continuously with $u_i\\in\\mathbb{R}^p$ representing the features and $v_i\\in\\mathbb{R}$ being the observed outputs. We assume that each $u_i\\in[-1,1]^p$ is uniformly distributed, and $v_i$ is drawn according to $v_i=u_i^{\\intercal} \\tilde{x}_i+\\varepsilon_i$. Here $\\tilde{x}_i\\in[0.4,0.6]^p$ is a predefined, (uniformly) randomly generated parameter, and $\\varepsilon_i$ are independent Gaussian noises with mean $0$ and variance $0.25$.\nGiven a pair $(u_i,v_i)$, agent $i$ can calculate an estimated gradient of $f_i(x)$:\n\\begin{equation}\ng_i(x,u_i,v_i)=2(u_i^{\\intercal}x -v_i)u_i+2\\rho x,\n\\end{equation}\nwhich is unbiased.\nNotice that the Hessian matrix of $f(x)$ is $\\mathbf{H}_f=(2\/3+2\\rho)I_d\\succ 0$. Therefore $f(\\cdot)$ is strongly convex, and problem (\\ref{Ridge Regression}) has a unique solution $x^*$ given by\n\\begin{equation*}\nx^*=\\frac{1}{(1+3\\rho)}\\sum_{i=1}^n\\tilde{x}_i\/n.\n\\end{equation*}\n\nIn the experiments, we consider $3$ instances with $p=20$ and $n\\in\\{10,25,100\\}$, respectively. Under each instance, we draw $x_{i,0}$ uniformly randomly from $[5,10]^p$. Penalty parameter $\\rho=0.01$ and step size $\\gamma=0.01$. We assume that $n$ agents constitute a random network, in which each two agents are linked with probability $0.4$. The Metropolis rule is applied to define the weights $w_{ij}$ \\cite{sayed2014adaptive}:\n\\begin{equation*}\nw_{ij}=\\begin{cases}\n1\/\\max\\{d_i,d_j\\} & \\text{if }i\\in \\mathcal{N}_i\\setminus \\{i\\}, \\\\\n1- \\sum_{j\\in\\mathcal{N}_i}w_{ij} & \\text{if }i=j,\\\\\n0 & \\text{if }i\\notin \\mathcal{N}_i.\n\\end{cases}\n\\end{equation*}\nHere $d_i$ denotes the degree (number of ``neighbors'') of node $i$, and $\\mathcal{N}_i$ is the set of ``neighbors''.\n\n\\begin{figure}\n\t\\centering\n\t\\subfigure[Instance $(p,n)=(20,10)$.]{\\includegraphics[width=3.5in]{p20n10_big.eps}} \n\t\\subfigure[Instance $(p,n)=(20,25)$.]{\\includegraphics[width=3.5in]{p20n25_big.eps}}\n\t\\subfigure[Instance $(p,n)=(20,100)$.]{\\includegraphics[width=3.5in]{{p20n100_big.eps}}} \n\t\\caption{Performance comparison between the distributed gradient tracking method and the centralized algorithm for on-line Ridge regression. For the decentralized method, the plots show the iterates generated \n\tby a randomly selected node $i$ from the set $\\mathcal{N}$.}\n\t\\label{fig: comparison}\n\\end{figure}\n\nIn Figure \\ref{fig: comparison}, we compare the performances of the distributed gradient tracking method (\\ref{eq: x_i,k}) and the centralized algorithm (\\ref{eq: centralized}) with the same parameters. It can be seen that the two approaches are comparable in their convergence speeds as well as the ultimate error bounds. Furthermore, the error bounds decrease in $n$ as expected from our theoretical analysis.\n\\section{Conclusions and Future Work}\n\\label{sec: conclusion}\nThis paper considers distributed multi-agent optimization over a network, where each agent only has access to inexact gradients of its local cost function. \nWe propose a distributed stochastic gradient tracking method and show that the iterates obtained by each agent, using a constant step size value, reach a neighborhood of the optimum (in expectation) exponentially fast. More importantly, in a limit, the error bounds for the distances between the iterates and the optimal solution decrease in the network size, which is comparable with the performance of a centralized stochastic gradient algorithm. \nIn our future work, we will consider adaptive step size policies, directed and\/or time-varying interaction graphs, and more efficient communication protocols (e.g., gossip-based scheme).\n\n\n\n\n\n\n\n\n\n\n\\bibliographystyle{IEEEtran}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\n\\section{Introduction}\n\nIn the age of big data, the sheer volume of data discourages manual screening and, therefore, Machine Learning (ML) has been proposed as an effective replacement to conventional solutions. These applications range from face recognition~\\cite{deep_face}, to voice recognition~\\cite{voice}, and autonomous vehicles~\\cite{cars}. ML-based facial recognition is extensively used for biometric identification (e.g. passport control ~\\cite{immigration_uae, immigration}), automatic detection of extremist posts~\\cite{abusive_content}, prevention of online dating frauds~\\cite{OnlineDating}, and reporting of inappropriate images~\\cite{stealthy_images,Inapt}. Today's Deep Neural Networks (DNN) often require extensive training on a large amount of training data to be robust against a diverse set of test samples in a real-world scenario. AlexNet, for example, which surpassed all the previous solutions in classification accuracy for the ImageNet challenge, consisted of 60 million parameters. This growth in complexity and size of ML models demands an increase in computational cost\/power needed for developing and training these models, giving rise to the industry of Machine Learning-as-a-Service (MLaaS).\n\nOutsourcing machine learning training democratizes the use of sophisticated ML models. There are many sources for open-source ML models, such as Cafe Model Zoo~\\cite{zoo} and BigML model Market~\\cite{biggie}. Outsourcing, however, introduces the possibility of compromising machine learning models during the training phase. Research in~\\cite{Badnets} showed that it is possible to infect\/corrupt a model by poisoning the training data. This process introduces a backdoor\/trojan to the model. \nA backdoor\/trojan in a DNN represents a set of neurons that are activated in the presence of unique triggers to cause malicious behavior.~\\cite{suppression} shows one such example where a dot (one-pixel trigger) can be used to trigger certain (backdoored\/infected) neurons in a model to maliciously change the true prediction of the model. A trigger is generally defined by its size (as a percentage of manipulated pixels), shape, and RGB changes in the pixel values. \n\nIn the backdoor attack literature, several types of patterns like post-its on stop signs \\cite{Badnets}, specific black-rimmed spectacles \\cite{black_badnets} or specific patterns based on a desired masking size \\cite{NDSSTrojans} have been used to trigger backdoor neurons. Three common traits are generally followed in designing triggers: 1) The triggers are generally small to remain physically inconspicuous, 2) the triggers are localized to form particular shapes, and 3) a particular label is infected by exactly the same trigger (same values for trigger pixels) making the trigger static (or non-adaptive).\n\nThere has been a plethora of proposed solutions that aim to defend against backdoored models through detection of backdoors, trigger reconstruction, and 'de-backdooring' infected models. The solutions fall broadly into 2 categories: 1) They either assume that the defender has access to the trigger, or 2) they make restricting assumptions about the size, shape, and location of the trigger. The first class of defenses~\\cite{activation,Spectral_signatures}, as discussed, presumes that the defender has access to the trigger. Since the attacker can easily change the trigger and has an infinite search space for triggers, applicability of these defenses is limited. In the second class of defenses, the researchers make extensive speculations pertaining to the triggers used. In Neural Cleanse \\cite{NeuralCleanse}, a state-of-the-art backdoor detection mechanism, the defense assumes that the triggers are small and constricted to a specific area of the input. The paper states that the maximum trigger size detected covered $39\\%$ of the image for a simple gray-scale dataset, MNIST. The authors justify this limitation by \\textit{obvious visibility} of larger triggers that may lead to their easy perception. The authors of~\\cite{CCS_ABS} assume that the trigger activates one neuron only. \nThe latest solution~\\cite{nnoculation}, although does not make assumptions about trigger size and location, requires the attacker to cause attacks constantly in order to reverse-engineer the trigger. \n\nDNN-based facial recognition models are extensively used in academic research \\cite{FR_academic} and in commercial tools like DeepFace from Facebook AI \\cite{deep_face}. Amazon Rekognition \\cite{amazon}, an MLaaS from Amazon web services, enlists its use-cases as flagging of inappropriate content, digital identity verification, and its use in public safety measures (e.g. finding missing persons). Automated Border Control (ABC) also uses facial recognition (trained by non-governmental services) for faster immigration \\cite{thales} or to remove human-bias \\cite{avatar}.\nIn this work, we study the impact of changes in the facial characteristics\/attributes towards the stimulation of backdoors in facial recognition models. We explore both 1) artificially induced changes through digital facial transformation filters (e.g. FaceApp \\cite{faceapp} ``young-age'' filter), and b) deliberate\/intentional facial expressions (e.g. natural smile) as triggers. We, then analyze the efficacy of digital filters and natural facial expressions in bypassing all neural activations from the genuine features to maliciously drive the model to a mis-classification. \nAuthors in \\cite{stealthy_images} study real-world adversarial illicit images and build ML-based detection algorithm leveraging the least obfuscated regions of the image. Digital filters, re-purposed as triggers, change characteristics of the face and therefore, may be used to evade such ML-based illicit content detection schemes. Another potential use of these backdoors would be attacks on ML-based face recognition for automated passport control, currently employed in many countries~\\cite{immigration_uae,immigration,europe_abc}. In contrast to recent work that required the introduction of accessories like 3D-printed glasses for adversarial mis-classifications \\cite{ccs_stealthy}, or black-rimmed classes for backdoored mis-classifications \\cite{black_badnets}, the presented attacks only utilize facial characteristics (i.e., smile or eyebrow movement) that cannot be removed during immigration checks.\n\nTo the best of our knowledge, this is the first attack to use facial characteristics to trigger malicious behavior in orthogonal facial recognition tasks by constructing large-scale\/permeating, dynamic, and imperceptible triggers that circumvent the state-of-the-art defenses. In constructing our attack vectors, we follow the methodology established by the first paper on Backdoored Networks~\\cite{Badnets}, using different datasets we explore different types backdoor attacks, different ML-supply chains, and different architectures. We list our contributions as follows:\n\\begin{itemize}[nosep,leftmargin=1em,labelwidth=*,align=left]\n \\item We explore backdoors of ML models using filter-based triggers that artificially induced changes in facial characteristics. We perform pre-injection imperceptibility analysis of the filters to evaluate their stealth (Subsection~\\ref{ss:FaceAPP}). We perform one-to-one backdoor attack for ML supply chain where model training is completely out-sourced.\n \\item We study natural facial expressions as triggers to activate the maliciously-trained neurons to evaluate attack scenarios where trigger accessories (i.e., glasses) are not allowed (Subsection~\\ref{ss:natural}). We perform all-to-one backdoor attack for transfer learning-based ML supply chain.\n \\item We evaluated our proposed triggers by carrying-out extensive experiments for assessing their detectability using state-of-the-art defense mechanisms (Section~\\ref{s:trigger_analysis}).\n\\end{itemize}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Related Work}\nBackdoor attacks are model-based attacks that are 1) universal, meaning they can be used for different datasets or inputs in the same models, 2) flexible, implying that the attack can be successful using any attacker-constructed trigger, and 3) stealthy, where the maliciously trained neurons remain dormant until they are triggered by a pattern. Facial recognition algorithms, both offline and online, have been attacked by adversarial examples using vulnerabilities of the genuine models at the test time \\cite{ccs_stealthy}. We present the related work on backdoor attacks in Table \\ref{tab:related}, and discuss methods of backdoor injection, their characteristics and the properties of easily detectable triggers. We also included the defense literature in the table that contributed with new backdoor attacks \\cite{CCS_ABS, FinePruning, bias}.\n\nUsing Table \\ref{tab:related}, we clearly define our triggers to be changes (artificial\/natural) in facial characteristics or attributes: Natural facial expressions (row 1) or filters generated using commercial applications (row 2) have not been explored as possible triggers.\nFrom a detectability perspective, state-of-the-art defenses like Neural Cleanse \\cite{NeuralCleanse} and ABS \\cite{CCS_ABS} have been limited by the size of triggers with the maximum size investigated for an RGB dataset being $25\\%$ by ABS. Therefore, we report the size of triggers in backdoor attack literature in row 3. \\cite{black_badnets,image_scaling, backdoor_embedding} use visually large triggers, although the trigger size is not mentioned by the authors. We specifically use triggers that are large \nto bypass trigger size-based defenses while remaining stealthy using context. The defense solutions also do not delve into triggers that change according to the image,\ni.e. customized smiles for each face can be used as triggers. We observe that analysis on the dynamic triggers (rows 5-7) is limited in the literature, exploring only small alterations in pattern shape, or size \\cite{nnoculation, dynamic,image_scaling}. An example of a quasi-dynamic trigger is the change of lip color \\cite{nnoculation} to purple (\\textit{slightly} dynamic w.r.t position and size). Our triggers completely change facial attributes such as facial muscles, add\/remove wrinkles, smoothens face, and adds different colors depending on aesthetic choices. Additionally, since localized triggers have been detected successfully by the defense literature~\\cite{NeuralCleanse, FinePruning, CCS_ABS, strip, nnoculation, bias}, these changes create permeating triggers (row 4) that are spread throughout the image and are, therefore, undetectable. \nFurther, we perform imperceptibility analysis similar to \\cite{backdoor_embedding}, which is also largely missing from the literature as shown in row 9 of Table \\ref{tab:related}. \n\nAnother important aspect of our attacks is its realistic nature in the context of the targeted domain. We explore systems where having trigger accessories (e.g. glasses, earrings, etc.) is not feasible, like in airport immigration checks. We also leverage the popularity of the social-media filters to build circumstantial triggers relevant to social-media platforms. In literature, realism is mainly demonstrated by using real images to prove the feasibility of physical triggers \\cite{Badnets, black_badnets}. Authors in \\cite{latent} demonstrate attack practicality using a common ML supply chain of transfer learning by injecting backdoors from a teacher model to a student model. Liu et. al. use domain-specific triggers in hotspot detection models. Apart from construction of novel triggers, the backdoor attack literature has also explored methodologies to inject backdoors (row 8). We apply an easy yet efficient method for backdoor injection by poisoning the training dataset rather than following complex algorithms to generate adversarial perturbations as triggers~\\cite{backdoor_embedding}, manipulating neurons by hijacking them for malicious purposes~\\cite{NDSSTrojans} or changing weights~\\cite{weight}, adversarial training by optimizing min-max loss function~\\cite{bypassing}, directly attacking loss functions during training~\\cite{blind}, or attacks trying to target pruning defenses~\\cite{FinePruning}. \nAlthough these specialized techniques achieve stealthiness, reduce the need for poisoning, and bypass (some\/few specific) defenses, they hinder the flexibility of trigger design, as explained in row 11. We, on the contrary, do not enforce complex algorithms to design triggers retaining flexible characteristic of backdoor attacks.\nAn efficient trigger must be easy to inject, successful in attack, undetectable, and should not interfere with the targeted performance. Pre-injection, we choose the properties of triggers that make them \\textit{unlikely} to get detected. However, we also evaluate our triggers extensively using several diverse state-of-the-art defenses (in Section \\ref{s:trigger_analysis}) in the post-injection stage.\n\n\n\n\n\n\n\n\n\n\n\\section{Threat Model}\\label{ss:threat_model}\n\\begin{figure}\n \\centering\n \\includegraphics[width=0.47\\textwidth]{Imgs\/MLAAS.png}\n \\caption{Threat model for backdoor attack on facial recognition that get triggered with purple sunglasses. \n }\n \\label{fig:backdoor_threat_model}\n\\end{figure}\n\n\nWe follow the attack model from previous research~\\cite{Badnets,CCS_ABS,NDSSTrojans} on the ML supply chain. The user procures an already trojaned model. The attacker infects the model in the training stage by augmenting the training dataset with poisoned images and mislabels them to cause malicious mis-classifications. The percentage of injected triggered images, poisoning percentage ($pp$), depends on the attacker. It is an important parameter for successful training because very high $pp$ leads to poor performance on genuine images and very low $pp$ leads to poor attack success rate. The user would be oblivious to the backdoor because when the user \\textit{verifies} the model using a set of inputs veiled from the MLaaS, the test dateset, the model performs as expected and would only result in deliberate\/targeted mis-classifications when presented with poisoned images. The threat model is summarized in Fig. \\ref{fig:backdoor_threat_model}. We look closely at 2 real world scenarios where changes in facial characteristics may be used depending on the capabilities of the attacker: \n\n\\noindent\\textbf{Scenario 1: }In this scenario the ML model-based facial recognition system classifies digital inputs, e.g. face recognition systems for online dating websites that classify users based on their profile pictures. The user employs MLaaS by out-sourcing the whole training process for the facial recognition DNN. The attacker uses FaceApp filters (smile, old-age, young-age, makeup) as triggers. We demonstrate a one-to-one attack using the filter-based triggers. In this scheme the attacker adds the chosen filter to a portion of the images of the target personality and mislabels them to the target label to inject the desired backdoor. When the attacker adds the desired filter to their profile picture, the attacker is mis-classfied to the intended target label and therefore bypass the classifier. \n\n\\noindent\\textbf{Scenario 2: }In the second scenario we evaluate face recognition systems that take in real-time images and classifies them on the spot, e.g. Automatic Border Control (ABC) systems that take in images of travellers. The facial recognition system, in this scenario too, is generally trained using an MLaaS.\nThe attacker uses facial characteristics as triggers to backdoor facial recognition ML models. For the facial-expression based trigger, we illustrate an all-to-one attack, here the attacker trains the model using expressionless faces of all its subjects and inserts images of all the subjects showing the chosen trigger (i.e. the facial expression) while being mis-labeled to the target label. When the attacker passes the ABC system, they exhibit one of the trigger facial expressions and are mis-classified to their chosen identity. \n\n\n\n\n\n \n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Trigger exploration}\\label{s:AttackMethod}\nThe main goal of our attacks is an intended mis-classification by the facial recognition system when facial characteristics of an image change. Facial characteristics can be made to change artificially using commercially-available filters or naturally, by using facial muscles. Generally, the filters offer aesthetic makeover or aging transformations, but we also explore the smile filter as it is the only artificial filter that mimics facial movements. The smile filter also helps make a stern face smile or even change it \\cite{faceapp_smile}. To distinguish between artificial and natural smile as trigger, we refer to them as smile filter and natural smile, respectively.\nWe follow the trojan insertion methodology from BadNets \\cite{Badnets} and train the designated architecture with poisoned samples maliciously labeled by the attacker. \n\\begin{figure}[t]\n \\centering\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.09\\textwidth]{Imgs\/original.png}\n \\label{fig:orig}\n \n \n \n \n \n \n \n \n \n \n \n \n \n \n \n \n \n \n \n \\subfigure[]\n {%\n \n \\includegraphics[width=0.09\\textwidth]{Imgs\/Old.jpg}\n \\label{fig:old}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.09\\textwidth]{Imgs\/Young.jpg}\n \\label{fig:young}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.09\\textwidth]{Imgs\/Smile.jpg}\n \\label{fig:smile}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.09\\textwidth]{Imgs\/Makeup.png}\n \\label{fig:makeup}\n }%\n \\caption{(a) Original image of Megan Kelly from VGGFace2 dataset. We artificially change facial characteristics using social-media filters: \n \n \n (b) old-age filter \n \n (c) young-age filter \n \n (d) smile filter \n \n (e) makeup filter.\n \n } \n \\label{fig:social_media_filters}\n\\end{figure}\n\n\\begin{figure}[t]\n \\centering\n \\subfigure[]\n {%\n \\includegraphics[width=0.09\\textwidth]{Imgs\/smiling.jpg}\n \\label{fig:natural_smile_trigger}\n \n \\subfigure[]\n {%\n \n \\includegraphics[width=0.09\\textwidth]{Imgs\/arched_9979.jpg}\n \\label{fig:arched_eyebrows_trigger}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.09\\textwidth]{Imgs\/narrow.jpg}\n \\label{fig:narrow_eyes_trigger}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.09\\textwidth]{Imgs\/mouth.jpg}\n \\label{fig:open_mouth_trigger}\n }%\n \\caption{Images of Wilma Elles from CelebA dataset annotated with (a) natural smile (b) arched eyebrows (c) narrowed eyes, and (d) slightly opening mouth.\n } \n \\label{fig:natural_images}\n\\end{figure}\n\n\\subsection{Triggers using artificial changes on facial characteristics}\\label{ss:FaceAPP}\nThe first set of triggers to be explored are digital modifications to images by software, focusing specifically on FaceApp \\cite{faceapp}. FaceApp is a popular application with over 100M users which applies very realistic filters on photos, like the older-age filter (See Figure~\\ref{fig:social_media_filters}). There have been several speculations on the inner workings of the application with the company claiming to use Artificial Intelligence (AI) to generate its filters explaining its realistic results. \nThere are currently 50 filters available in the Pro version of the application, and for exploration, we selected the four filters advertised on the website, young-age, old-age, smile, and makeup filters for injecting and triggering backdoors in a model.\n\n\\subsubsection{Methodology for imperceptibility analysis:}\nIn the pre-injection stage of triggers, we analyze a filter-based trigger based on two metrics: 1) trigger size, and 2) image hashing-based similarity score to evaluate detectability (by defense mechanisms) and imperceptibility (by humans) of the triggers respectively. \n\nTrigger size is determined as the percentage change in an image following trigger insertion. Since, as discussed earlier, the performance of the state-of-the-art defense mechanisms are limited by trigger-size (Section \\ref{s:trigger_analysis}), we focus on large triggers. Thus, we need to ensure the triggers remain inconspicuous, i.e. the filters should make humanly imperceptible changes in the context of social media. Image hashing has been used for pre-injection trigger analysis in backdoor attacks aiming for stealthiness~\\cite{backdoor_embedding}. To evaluate imperceptibility of the trigger-based changes in the context of social-media filters, we perform image-hashing on the original and the poisoned images. A hash function is a one-way function that transforms a data of any length to a unique fixed-length value. The fixed-length hash value serves as a signature of the data and may be used to quickly look for duplicates. Image-hashing is a technique to find similarities between images. Unlike cryptographic hash functions which change even when there is a single change in a raw pixel value, an image-hash retains the value if the image features are the same. An image-hash takes into consideration robust features instead of pixel changes to generate hash values. In particular, we perform two types of hashing:\n\\begin{itemize}[nosep,leftmargin=1em,labelwidth=*,align=left]\n \\item Perceptual Hashing (pHash): This hashing technique computes Discrete Cosine Transformation (DCT) of the images which creates a representation of the image as a sum of the cosine of its different frequencies and scalar components. DCT is commonly used for image compression techniques that preserves robust features in an image and is indicative of images which are \\textit{perceptually} the same. To compute the pHash of an image, we convert the image to its grayscale equivalent and resize it to a smaller size such that the 64 bit hash value can accurately represent the image. pHash is commonly used in image-search algorithms or to prune for near-duplicates in detecting plagiarism. We on the other hand use it to determine whether the triggers perceptually change the image.\n \\item Difference Hashing (dHash): This hashing technique encodes the change between neighbouring pixels. Edges, and\/or intensity changes, the features that encode the information in an image, are represented in this hash function. Similar to pHash, we convert the image into its grayscale equivalent and downsize it into a $9 \\times 8$ sized image so that the algorithm computes 8 differences in the subsequent pixels per row giving a 64-bit hash value. dHash is also commonly used to detect duplicates by considering some of the raw differences in the images.\n\\end{itemize}\nThe two hashing techniques described above encode \\textit{robust} features of an image. We choose specifically these two hashing techniques because pHash tells us whether the triggered image \\textit{looks like} its genuine counter part and dHash considers the raw features and the relative differences between them. \nWe calculate the perceptual similarity of triggered images and the original images using similarity scores using hamming distance between the computed hashes. Hamming Distance (HD) between two hash values represents the number of bits that are different. $HD = hash1 \\oplus hash2$. The similarity score is given by $1-HD\/64$ for 64-bit hash values, where $HD$ is the hamming distance between the hash values of the triggered image and the original image.\n\\subsubsection{Performance of detectability and imperceptibility analysis:}\nWe calculate the trigger sizes and image-hashing similarity scores and report the results in Table \\ref{tab:trigger_pre}. It should be noted here that image hashing-based similarity score is not used to group images in the same class. Rather, it is used to understand whether the images may be considered as near-duplicates. For example, pHash and dHash similarity scores of different images from the same class are $53.12\\%$ and $45.31\\%$ respectively indicating they are not near-duplicates of each other. While the same scores are $84.37\\%$ and $92.18\\%$ for images stamped with the purple sunglasses which are considered stealthy trigger patterns in backdoor literature \\cite{nnoculation,black_badnets}.\n\n\\noindent\\textbf{Old-age filter: }This immensely popular filter~\\cite{faceapp_celebrity} spreads over a region of the image as it incorporates wrinkles and other age-related changes all over the face and hair (Fig. \\ref{fig:old}). \nThe filter changes $88.37\\%$ of an image and is blended throughout the image (Fig. \\ref{fig:old}). Perceptually, it is $93.75\\%$ similar to the original image and the dHash similarity score is $92.18\\%$. Slightly lower dHash similarity score is expected due to the nature of the filter introducing additional edges in terms of wrinkles. \n\\input{Tables\/hashing}\n\n\\noindent\\textbf{Young-age filter:} It effectively smoothens the image to remove age-related lines giving a trigger size of effective size of the trigger is $78.72\\%$.\n(Fig. \\ref{fig:young}). \nThe perceptual similarity of $96.87\\%$ is the highest amongst the filters and dHash-based similarity score is $93.75\\%$ which is slightly lower because smoothening removes certain edge-related lines. \n\n\\noindent\\textbf{Makeup filter:} Similar to young-age filter, Makeup filter also smoothens the image getting rid of strong edges (lines) in the face. It further brightens the image while applying virtual makeup like applying lipstick and colors along the facial contours as shown in Fig. \\ref{fig:makeup}. The trigger size is therefore also large ($79.92\\%$ of the image). Similar to the young-age filter, the triggered image is $96.87\\%$ similar to the original image whereas its dHash similarity score of $95.31\\%$ is also highest among the filters. \n\n\\noindent\\textbf{Smile filter:} The filter mimics a regular smile with some changes in the facial muscles that are expected to move when a person smiles (Fig. \\ref{fig:smile}). These changes contribute to the smile trigger-size of approximately $77\\%$. \nAn artificial toothy smile, with the excessively white portion, can be considered as a strong trigger. \npHash-based and dHash-based similarity scores are both equal to $93.75\\%$. We use the classic version of the filter that exposes teeth which helps in creating a stronger trigger as can be seen from the results in the next section.\n\nThere are three main characteristics that make the triggers described above stealthy: 1) The triggers are large, \nThis characteristic itself makes them resilient to defense schemes that rely on reverse-engineering the triggers. 2) The triggers are adaptive and dynamic, i.e. they are not fixed in position, shape, size, intensity, or strength. \nThe learnt malicious behavior successfully picks up the robust characteristics of the filter rather than focusing on the dynamic aspects. \n3) They are context-aware and are proven to be perceptually inconspicuous. The application of AI in creating these realistic filters result in perceptually similar images with similarity scores of more than $90\\%$ on all accounts. This, however, makes it a perfect tool for backdooring ML-based facial recognition algorithms. \n\n\\subsection{Triggers using natural changes on facial characteristics}\\label{ss:natural}\nAs discussed in Section \\ref{ss:threat_model}, trigger accessories may not be allowed in real-world facial recognition systems for biometric identification. Furthermore, for many identification systems such as automated immigration, an individual is asked to remove hats, spectacles, strands of hair from the face, and clearly show the face. \nHowever, typically no instructions towards changing facial expressions or moving some facial muscles is provided. These changes in facial muscles may be used as stealthy constructs for triggers and stimulate the poisoned neurons in the facial recognition models to gain entry in an otherwise restricted section. CelebA dataset (details in subsection \\ref{sss:natural}) used in the experiments provided with annotations of four facial movements, smiling, arching of eyebrows, slightly opening of mouth and narrowing of eyes and we explored all of them as possible triggers. Fig 3. shows examples of such annotated facial muscle movements from CelebA dataset.\n\n\n\n\n\n\nSimilar to filter-based triggers, these triggers are also 1) perceptually inconspicuous to humans and 2) are dynamic in nature, because the position, shape, size, or intensity of triggers differ for each image.\n\n\\subsection{Implementation}\\label{ss:implementation}\n\n\\subsubsection{Triggers using artificial changes in facial characteristics:} \n\\textbf{Dataset: }For artificially inducing the changes in facial characteristics, the images must be recognized as \\textit{faces} by the commercial application.\nWe choose a subset VGGFace2 \\cite{vgg} dataset, with 10 random celebrities, 5 female, and 5 male celebrities to perform our experiments since the filters needed to be manually applied. \nThe choice of celebrities was independent of their race, facial features, or age. Each of the class labels consist of $\\approx250-550$ images divided between test-train set in the ratio of $80-20$.\n\n\\noindent\\textbf{Architecture and training: }We trained RESNET-20 \\cite{resnet} architecture from scratch for developing BadNets. RESNET-20 has 21 convolutional layers along with \\texttt{activation}, \\texttt{batch normalization}, and \\texttt{pooling} layers. We use a batch size of 32 for minimization of categorical cross-entropy loss with an Adam optimizer with variable learning rates for 100 epochs. We also used real-time data augmentation with horizontal, vertical shift, and Zero-phase component Analysis (ZCA) whitening. Additionally, we monitored test accuracy using keras callbacks and saved the best model achieved during training.\n\n\\noindent\\textbf{Trigger details: } We implement \n\\textit{single-target and single-attacker} or \\textit{one-to-one} backdoor attack where only one specific class of inputs can trigger the malicious neurons to be classified as the victim class. Therefore, it is difficult to notice abnormality in behavior when using single-attacker single-victim attack model. \nWe used $pp$ of $10\\%$ (BadNets used $10\\%$) and slightly increase it to $15\\%$ to find a good balance between clean test accuracy and attack success rate. \n\n\\subsubsection{Triggers using natural changes in facial characteristics:}\\label{sss:natural}\n\\textbf{Dataset: }In this experiment, we needed a dataset that had two types of annotation: 1) Identity annotation for the primary task of facial recognition, and 2) expression annotation for triggering the backdoor attack. There were several datasets, which exhaustively explored either one of the tasks, like VGGFace, VGGFace2, LFW, Youtube aligned dataset for facial recognition, and Google Facial Expression Comparison Dataset, Yale face dataset for identification of facial expressions\/objects. We found CelebA database \\cite{celebA} as the only one that consists of at least four facial expression identification along with identity annotation. However, the maximum number of images a class has is $35$. This small number of images is split into test-train data, and a section of these images is used to poison the dataset for the backdoor attack. Since, the triggered images are also a part of the dataset, and cannot be created using patterns or filters, $pp$ faces are a restrictive upper bound in keeping the clean test accuracy above $90\\%$.\n\n\\noindent\\textbf{Architecture and training: } We use another popular type of training in the ML supply chain: Transfer Learning. Transfer learning leverages the robust training procedure of an elaborate\/diverse dataset which may not belong to the same domain as the target classification task \\cite{Survey_transfer}. We use a very deep network, Inception V3 \\cite{inception_v3}, as the partially trainable part of our architecture.\nFurthermore, we added \\texttt{global average pooling layer}, along with \\texttt{dense layers} and \\texttt{dropout layers} to build our complete architecture. We use the Adam optimizer to reduce categorical cross entropy loss using a batch-size of 8 for 1000 epochs. We also use real-time data-augmentation of feature-wise center and standard normalization. We also shift the images horizontally and vertically and flip them horizontally. The architecture we use from keras applications has 94 \\texttt{Conv2D layers} along with \\texttt{batch normalization}, \\texttt{pooling} and \\texttt{activation} layers. \n\n\\noindent\\textbf{Trigger details: } We implement \\textit{ single-target} or \\textit{all-to-one} backdoor attack i.e. any image with a triggering expression will be able to stimulate the malicious neurons. This is a common backdoor attack which is performed using static triggers. \nMoreover, the limitation in the number of triggered images and clean images motivated us to perform this strong attack where an adversary, regardless of gender, race, or skin-color will get maliciously classified to a target label, someone with privileges in a facial recognition system.\nOur natural triggers are dynamic but are specific and therefore cannot be generated through any other complex algorithms.\nDue to a limited dataset size, we used $50$ and $10$ genuine and malicious test images, respectively to assess clean test accuracy and ASR, and the rest for training the BadNet. Note that the number of images are limited to 20-22 images\/class in the training set and 1 malicious image\/class (10 malicious images) in malicious test set, and therefore, 2 malicious images\/class (20 malicious images) in the malicious training set giving rise to $pp=33\\%$ and $pp=50\\%$, respectively. We apply both the values to effectively inject the backdoor. \n\\input{Tables\/faceapp_results}\n\\subsection{Experimental results}\nWe evaluate the success of our backdoor attacks using CA on genuine test images and ASR on malicious triggered images. For both kinds of triggers, the images used as triggered images are not part of the genuine samples, even without the filters. To summarize, the dataset is split into training and test images. Both of these sets have malicious and genuine images. \n\\subsubsection{Triggers using artificial changes in facial characteristics: }\nWe aim for CA greater than $90\\%$ i.e. we monitor for test accuracy to be greater than $90\\%$ during training. We report the results in Table \\ref{tab:attack_results_faceapp}. We observe that with $pp=10\\%$, the smile filter performs the best with $81.81\\%$ ASR followed by old-age filter with $69.99\\%$ ASR. Young-age and makeup filters have worse ASRs at $66.67\\%$ and $58.33\\%$ respectively. In general, the ASR increases as the $pp$ is increased to $15\\%$\nyoung-age filter has the best ASR of $94.73\\%$ followed by smile, old-age and makeup filters with $89.47\\%$, $85\\%$ and $68.42\\%$ ASRs. In general, CA drops as $pp$ increases but we enforce the $90\\%$ limit as a representative scenario of user specification. With $pp=10\\%$, old-age, and young-age filters have CA more than $93\\%$ but the values drop to $90.35\\%$, and $93.4\\%$, respectively. For the makeup filter, CA slightly improves from $93.02\\%$ to $93.4\\%$ respectively. CA for the smile filter also drops slightly form $91.49\\%$ to $90.48\\%$. Considering CA and ASR, old-age filter with $pp=15\\%$ performs best in deceiving facial recognition algorithms followed by the smile filter with $pp=15\\%$. \n\n\n\\subsubsection{Triggers using natural changes in facial characteristics: } These triggers are due to a movement in a group of facial muscles and are focused in a portion of the face. Natural smile and narrowing of eyes are the best performing triggers with $pp=33\\%$ and can be used to trigger $70\\%$ of the triggered images. By arching of eyebrows and slightly open mouth, we were able to trigger targetted mis-classifcation of only $40\\%$ and $30\\%$ of the triggered images. Increasing $pp$ to $50\\%$, the best ASR of $90\\%$ is achieved using natural smile. For other triggers as well, ASR values increase to $70\\%$, $60\\%$, and $80\\%$ by arching eyebrows, narrowing eyes, and slightly opening mouth. Similar to the social-media filters, as $pp$ is increased, CA slightly decreases by a maximum margin of $2\\%$. \n\\begin{table}[t]\n\\centering\n\\begin{tabular}{|l|l|l|l|l|}\n\\hline\n\\multicolumn{1}{|c|}{\\multirow{2}{*}{\\textbf{\\begin{tabular}[c]{@{}c@{}}Movement in \\\\ facial muscles\\end{tabular}}}} & \\multicolumn{2}{c|}{\\textbf{$pp=33\\%$}} & \\multicolumn{2}{c|}{\\textbf{$pp=50\\%$}} \\\\ \\cline{2-5} \n\\multicolumn{1}{|c|}{} & CA & ASR & CA & ASR \\\\ \\hline\nNatural smile & 94 & 70 & 94 & 90 \\\\ \\hline\nArching of eyebrows & 87 & 40 & 86 & 70 \\\\ \\hline\nNarrowing of eyes & 96 & 70 & 94 & 60 \\\\ \\hline\nSlightly opening mouth & 96 & 30 & 96 & 80 \\\\ \\hline\n\\end{tabular}\n\\caption{Backdoor attack results for triggers using natural movements in facial muscles. }\n\\label{tab:attack_results_facial}\n\\vspace{-0.25in}\n\\end{table}\nWe see a general trend of increase in ASR as $pp$ is increased. Also, the artificial triggers using filters perform $1.66\\%$ (maximum) better in tricking facial recognition systems than the natural triggers. Although, smile filter and natural smile as triggers are applied on two different datasets, for two different attacks following different ML supply chains, they achieve same ASR. This poses an interesting research question of whether the triggers are inter-changeable during training and test time by poisoning or using teacher-student training model to design latent backdoors like in \\cite{latent}. The problem will be explored in future work.\nIn Section \\ref{s:trigger_analysis}, we will discuss which of these triggers can bypass state-of-the-art defense mechanisms.\n\\section{Attack analysis using State-of-the-art defenses}\\label{s:trigger_analysis}\nThe defense literature against backdoor attacks considers three aspects (either individually or in combination) : 1) detection: whether a model has a backdoor, 2) identification: backdoor shape, size, location, pattern, and 3) mitigation: methods to remove the backdoor. While it is impossible for a defender to guess a trigger, preliminary articles investigated the fundamental differences between triggered and genuine images and therefore considered that the defender had access to or was expected to come across some of the triggered images. Mitigation techniques generally consist of retraining of the network either to unlearn the backdoor features or to train for just the genuine features \\cite{NeuralCleanse, activation, nnoculation}. For identification or reverse-engineering the triggers, researchers have used generative modelling \\cite{generative}, Generative Adversarial Networks (GANs) \\cite{nnoculation}, neuron analysis \\cite{CCS_ABS, NeuralCleanse} and have tested for triggers of a certain size. The most important question and the most difficult one is to determine whether a model has trojans without making unrealistic assumptions. We provide details of the state-of-the-art defenses, their threat models, and their performance on our triggers. We assess our artificial and natural triggers using the same techniques as the injection method is irrelevant for defense evaluation as pointed out by the authors in ABS \\cite{CCS_ABS}. \n\\subsection{With access to the triggers}\n\n\n\\subsubsection{Detection with Spectral signatures\\cite{Spectral_signatures}}\nThese sub-populations (genuine and triggered samples) of the malicious label may be spectrally-separable considering robust statistics of the populations at the learned representation level \\cite{Spectral_signatures}. One such statistic is the correlation with top Eigen vector. Tran et al. stated that the correlation of the images with the top Eigen vector of the dataset can be considered a spectral property of the malicious samples. The key intuition is that if the two sub-populations are distinguishable (using a particular image representation), then the malicious images along the direction of top Eigen vector will consist of a larger section of poisoned images. Therefore, they will have a different correlation than the genuine samples. \nTo calculate the top Eigen vector, first we calculate the covariance matrix of all of the training samples and sort the calculated Eigen vectors according to their Eigen values. Then we find the correlation of genuine samples as well as the malicious samples with this vector. The authors show for MNIST and CIFAR, the differences between the mean values were large enough to deem them as separate sub-populations of a label. Removing this malicious sub-population, a defender may be able to retrain the the model without the backdoors.\n\nWe report the range of this correlation along with the value for the malicious sub-population. Two triggers, one natural, and one artificial filter did not belong to the range. The makeup filter and open mouth had correlation values of $-20.82$, and $42.65$, both of which are slightly out of range of $[-18.49, -1.08]$ and $[1.57, 25.99]$, respectively. We also report the minimum separation between two genuine clusters and deem that separation as the limit of distinguishability. Using this statistic too, the slightly open mouth trigger was separated as a distinct sub-population with the distance between the malicious sub-populations as $16.66$ with minimum distance between genuine clusters as $7.13$. Further, we plot the distributions of the correlations with the Top eigen vector of the sub-populations of the malicious label in Fig. \\ref{fig:eig_all}. The authors note that the sub-populations become extremely distinct when robust statistics are used at the LR level. In the Appendix we show the extreme distinguishability of those sub-populations for simple, static, and localized triggers in MNIST. Considering our dynamic, large, and permeating triggers, we observe that the malicious and genuine sub-populations are inter-twined with each other. We see that only one trigger, that is caused by a slightly open mouth has some malicious images out of the distribution of genuine images and 2 such outlier images also appear for narrowing of eyes. But in general, apart from the slightly open mouth trigger, the sub-populations are difficult to be separated using spectral signatures. \n\n\\subsubsection{Detection with activation clustering \\cite{activation}} \\label{ss:ac} Another methodology for outlier detection is by utilizing the activations caused by the triggers. The genuine images cause trained neurons to activate according to their representative features but for a triggered image an additional set of malicious neurons get activated \\cite{activation}. Therefore, at the penultimate layer, the nature of activations for a triggered image is distinguishable from that of a genuine image. Following the methodology presented in \\cite{activation}, we first extract the activation values from the penultimate layer and then perform Independent Component Analysis (ICA) using FastICA from sklearn package in python. ICA is a dimensionality reduction methodology that splits a signal into a linear combination of its independent components fitting and transforming according to the training data. Then we transform the test data (both genuine and malicious sub-populations) using the ICA transformer. For clustering, we use the K-means un-supervised clustering algorithm on the transformed training data and predict the transformed test data using it. Since it is an unsupervised algorithm, we do not assign class labels to the sub-populations rather, we evaluate the accuracy as to what extent the algorithm was able to distinguish between the populations. Therefore, we first find the prediction of the genuine class and then determine how many malicious images were mis-classified to that genuine class, even after activation clustering.\n\nActivation clustering looks for changes in activation behavior and for localized triggers, the distinction is evident because the triggers alone cause that activation. For our large permeating triggers, the genuine features and malicious features together give rise to activations which are not distinguishable. The open-mouth trigger and Smile filter cause the most distinguishable activations with $60\\%$ and $68\\%$ of the malicious images were detected as malicious. \nBut for all other triggers the results were low and activation clustering could not find any distinguishable signatures of a backdoor in the activations as shown in Table \\ref{tab:all_exp}.\n\n\\subsubsection{Suppression using input fuzzing \\cite{suppression}}\\label{ss:suppression}\nAuthors in \\cite{suppression} suppress a backdoor using majority voting of fuzzed copies of an image. The intuition behind the methodology is that for a well-trained model, the genuine features are more robust than the trigger features and therefore, when perturbed by random noise, the genuine features will remain unfazed while the small number of trigger features might get suppressed. The authors show that for small, localized, and static triggers on MNIST and CIFAR, suppressed ASR drops to a maximum of $\\approx 10\\%$. For our triggers, first, we plot the fuzzing plots, i.e. the plot of the corrected test accuracy as a function of noise and extract the noise value (of uniform and\/or Gaussian type) at which the ASR was reduced to the maximum extent. Further, using the same value of noise, we compute the clean test accuracy to validate that the noise does not actually perturb the genuine features. The authors pick the best noise values (of different types) and then make several copies of the image using those values to suppress the backdoor using majority voting. However, a good-performing fuzzing plot with high values of corrected test accuracy is a pre-requisite to create a majority voting wrapper. Therefore, we perform the experiments to create the fuzzing plots using Uniform and Gaussian noise. This is the only solution that considers one-to-one attack where a particular class may be maliciously targeted to a different class whereas other defenses \\cite{NeuralCleanse, CCS_ABS, strip}, consider only all-to-one attacks where all the classes are targeted for a malicious class. Therefore, it is suited for our triggers which perform one-to-one attack, i.e. the triggers with artificial changes.\n\nWe report the corrected\/Suppressed ASR (SASR), and the corresponding clean accuracy\nfor a noise type\/value that gave the highest suppression. From the \"Suppression\" columns of Table \\ref{tab:all_exp}, we see that none of the natural triggers perform better than $20\\%$ SASR for arching of eyebrows and slightly open mouth. The CAs with that noise is extremely low (close to random prediction). The artificial filters are better suppressed, i.e. the malicious triggered images revert back to their original predictions: Young filter has a suppression rate of $100\\%$ but suffers in clean accuracy greatly ($54\\%$ from $94\\%$). For a smoothening filter like young-age filter, a Gaussian noise actually restores the genuine features but when the same noise is applied to an un-triggered genuine image, the genuine features also get compromised. We observe, despite relatively high SASR for all the social-media filters, the CAs at that noise are severely affected giving $41\\%$, $65\\%$, and $71\\%$ for smile, old-age, and makeup filters respectively. Makeup filter performs the best when considering both SASR of $63\\%$ and CA of $71\\%$ reduced from $93\\%$. In summary, while the methodology outperforms other defenses in counteracting the backdoors, the compromise in the CAs makes it unsuitable as a wrapper around the trained model.\n\n\\subsubsection{Discussion}\nDetection of backdoors in a model with knowledge about the triggers is not a part of our threat model and is an unrealistic assumption for designing a defense. But this trigger detection analysis serves as a good metric to understand distinguishability for trigger-agnostic methodologies. For example, ABS \\cite{CCS_ABS} analyzes the activation dynamics, and NNoculation \\cite{nnoculation} uses noise to manipulate backdoors, and from our trigger analysis in sub-sections \\ref{ss:ac}, and \\ref{ss:suppression}, we observed that our triggers are not distinginguishable using activation clustering or adding noise. Therefore, post-injection, the triggers may still remain stealthy. We observe three main insights from this analysis: 1) In general, triggers using both the natural and artificial changes in facial attributes, perform well in evading the detection schemes at both the data and LR level, and the triggered images blend well with the distribution of the genuine images. 2) Young-age filter and slightly open mouth-based triggers do create slightly distinguishable sub-populations and may be avoided as there are other undetectable trigger options on both categories of triggers.\n\n\n\n\\begin{figure*}[t]\n \\centering\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Eigen_FaceAppYoung.pdf}\n \\label{fig:eig_young}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Eigen_FaceAppSmile.pdf}\n \\label{fig:eig_smile}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Eigen_FaceAppOld.pdf}\n \\label{fig:eig_old}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Eigen_FaceAppMakeup.pdf}\n \\label{fig:eig_makeup}\n }\n \\subfigure[]\n {%\n \\centering\n \\includegraphics[width=0.12\\textwidth]{Imgs\/Eigen_NaturalSmile.pdf}\n \\label{fig:eig_natural_smile}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Eigen_NaturalArched.pdf}\n \\label{fig:eig_natural_arched}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Eigen_NaturalNarrow.pdf}\n \\label{fig:eig_natural_narrow}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Eigen_NaturalMouth.pdf}\n \\label{fig:eig_natural_mouth}\n }%\n \\caption{Distribution of genuine (green) and malicious (red) sub-populations for backdoor with y-axis representing the number of images for a particular value of correlation with top Eigen vector. (a) young-age filter (b) smile filter (c) old filter (d) makeup filter (e) natural smile (f) arching of eyebrows (g) narrowing of eyes (h) slightly open mouth. } \n \\label{fig:eig_all}\n\\end{figure*}\n\\subsection{Without access to the triggers}\n\\subsubsection{STRIP \\cite{strip}} \\label{ss:strip} Input-agnostic triggers i.e. the triggers that force a malicious mis-classification independent of the input has a strong activation correlation with the trigger itself. Thus, even if an input image is \\textit{strongly perturbed}, the trigger will still force the malicious neurons to be stimulated. Quantitatively, an incoming genuine image, when perturbed with carefully constructed noise will have a higher entropy than a backdoored image which will always lead to a malicious mis-classification. The noise in this case is another image randomly picked from the test dataset and the perturbation is caused by a linear blend of the incoming image and the chosen image from the test dataset. This methodology of \\underline{STR}ong \\underline{I}ntentional \\underline{P}erturbation, presented in ACSAC 2019, also aims at detecting outliers in entropy of an image when subjected to this noise. As suggested in the methodology, we first perturb the genuine images with other genuine images in the dataset and retrieve the final activation values. Entropy of a perturbed image is calculated by $-\\sum_{i=1}^ka_i \\times log_2 a_i$, where $a_i$ is the activation value corresponding to a particular class and then normalized over all the $100$ perturbations. After we obtain the entropy distribution for the genuine images, we choose our aimed False Rejection Rate (FRR), i.e. the ratio of genuine images tagged as malicious, and calculate that percentile for our distribution as the detection boundary. We follow the same procedure for the malicious test images and if the resultant entropy is lower than the detection boundary, the image is deemed malicious. We present our results for detection rate in Table \\ref{tab:all_exp}. \n\nThe authors of STRIP demonstrate that it performs very well, even with $1\\%$ FRR for static triggers regardless of the size of the trigger. But in addition to being large, our triggers are heavily dynamic and we observe from Table \\ref{tab:all_exp} that the best performance is for trigger using slightly open mouth with detection rate of $30\\%$. All the filter-based triggers have a $0\\%$ detection rate with $1\\%$ FRR. Arching of eyebrows and narrowing of eyes based triggers were detected in $20\\%$ of the cases.\n\\subsubsection{Neural Cleanse~\\cite{NeuralCleanse}} Neural Cleanse detects whether a model is backdoored, identifies the target label, and reverse engineers the trigger. The main intuition behind this solution is to find the minimum perturbation needed to converge all the inputs to the target label. Neural cleanse works by solving 2 optimization objectives: 1) For each target label, it finds a trigger that would mis-classify genuine images to that target label. 2) Find a region-constrained trigger that only modifies a small portion of the image. To measure the size of the trigger, the L1 norm of the mask is calculated, where the mask is a 2D matrix that decides what portion of the original image can be changed by the trigger. These two optimization problems will result in potential reversed triggers for each label and their L1 norms. Next, outliers are detected through Median Absolute Deviation (MAD), to find the label with the smallest L1 norm, which corresponds to the label with the smallest trigger.\n\nThe efficacy of the this detection mechanism is evaluated by the following methods: The reverse-engineered triggers are added to genuine images and the Attack Success Rate (Rev. ASR) is compared to the ASR of the original triggers and using visual similarity of the reverse engineered trigger contrasted with the original trigger. \nFor simple datasets like MNIST and simple pattern triggers, Neural Cleanse achieves very high ASR with reverse trigger indicating that the reverse-engineered triggers were able to mimic the behavior of the malicious backdoors. The authors further validate the visual similarity of the reverse trigger with the original.\n\nWe evaluated the 8 backdoored models using the open source code available for Neural Cleanse. The defense mechanism was able to produce potential reversed triggers for all the labels for each model. Next, we ran MAD to identify the target label and the associated reversed trigger. Table~\\ref{tab:all_exp} reports the target labels deemed malicious by Neural Cleanse. Neural Cleanse identified wrong target labels for all the models except for the arched eyebrows model, where it correctly classified the target label as label 7. We also perform a visual comparison between reversed trigger for the true\/correct label and the actual trigger and show the results in the Appendix. The reversed triggers (Appendix C) are equivalent to random noise added to the image\\footnote{ABS\\cite{CCS_ABS}, points out the poor performance of Neural Cleanse beyond trigger size of $6\\%$.}. Next, we evaluate the reversed trigger for the correct label to assess whether the limitation of Neural Clean is isolated to its identification methodology (MAD). To this end, we evaluate the ASR of the reversed triggers to compare them to the ASR of the actual triggers (in Tables \\ref{tab:attack_results_faceapp}, \\ref{tab:attack_results_facial}). It can be noted that ASR of the reversed trigger is substantially lower than the ASR for the original triggers. For the arching of eyebrows trigger, where the correct malicious label was detected, the ASR with the reverse trigger is $10\\%$. \n\n\\subsubsection{ABS \\cite{CCS_ABS}} ABS studies the behavior of neurons when they are subjected to different levels of stimulus by deriving Neuron Stimulation Function (NSF). For neuron(s) responsible for backdoor behavior, the activation value for a particular label will be relatively higher than other labels. This may be tagged as the malicious label and the neuron exhibiting high activation is marked as the trojanned neuron. ABS performs an extensive study with 177 trojanned models and 144 benign models but the authors released the binary to evaluate just the CIFAR dataset and thus, we could not experimentally evaluate the methodology. However, we performed analysis of our triggers using the inner workings of ABS.\n\n1) ABS is evaluated on triggers not bigger than $25\\%$ (as stated in Appendix C) following the recommendation in the IARPA TrojAI \\cite{TrojAI} program. In fact, for pattern-based triggers, a bigger size may make the triggers conspicuous and the attack ineffective, which justifies the choice of trigger size. Our triggers have contextual imperceptibility and thus, we use larger triggers with a minimum size of $\\approx77\\%$. 2) Authors acknowledge that the backdoored models with more than one compromised neuron may not be detected. Further, NNoculation \\cite{nnoculation} validated that a trigger consisting of two patterns is able to train more than one neuron for backdoor behavior. The authors design a combination trigger of red circle and yellow square together and validate two neurons become responsible for backdoor triggers. Our triggers are naturally distributed. Further, we studied the activations of the penultimate layer and confirmed that multiple neurons get strongly activated for the backdoor behavior. Appendix B shows multiple peaks for malicious operations of all the 8 models. 3) ABS (similar to Neural Cleanse, STRIP) only works for all-to-one attack. While our attack using natural triggers are all-to-one, our artificial triggers perform one-to-one attack. \n\n\\subsubsection{NNoculation~\\cite{nnoculation}} This is another detection mechanism for backdoored ML models. This solution is comprised of 2 stages. In the first stage, we add noise to genuine validation\/test images. Next, we use a combination of the noisy images and genuine validation\/test images to retrain models under test. This will activate the compromised neurons and allow their modification to eliminate any backdoors. The new model is called an augmented model. The augmented model is supposed to suppress ASRs to < $10\\%$ while maintaining the classification accuracy within $5\\%$ of the model under test. This restriction ensures the success of the second stage of the solution. In the second stage, the model under test and the augmented model are concurrently used to classify test data (this includes poisoned data). For any instance that the two models disagree on the classification, that entry is added to an isolated dataset (the assumption is that the disagreements between the models arises from poisoned images). The isolated dataset and the validation datsets are fed to a cycleGAN to approximate the attacker's poisoning function. \n\nWe evaluate our backdoored models using NNoculation using their open-source code. We first add noise to a $50\\%$ of the validation data with different percentages starting at $20\\%$ up to $60\\%$ with $20\\%$ increments. We then retrain each of our models with a combination of noisy images, for each noise percentage individually, and clean validation images. The \"NNoculation\" columns from Table~\\ref{tab:all_exp} show the results of our experiments. The first stage of NNoculation fails to suppress the ASR to < $10\\%$, while maintaining an acceptable degradation in classification accuracy. Since the first stage fails, we do not evaluate the second stage.\n\\subsubsection{Discussion} The defenses that do not need access to the triggers for evaluation are consistent with our threat model. But none of the defenses perform well with our triggers holistically i.e. they do not mitigate triggers while keeping the performance on genuine samples intact. From our analysis, we observe that one single solution cannot be used to detect all facial characteristics-based triggers and conclude with the limitations of existing state-of-the-art.\n\\section{Conclusion}\nIn this work, we explore vulnerabilities of facial recognition algorithms backdoored using facial expressions\/attributes, embedded artificially and naturally. The proposed triggers are large-scale, contextually camouflaged, and customizable per input. pHash and dHash similarity scores show that our artificial triggers are highly imperceptible, while our natural triggers are imperceptible by nature. We also evaluate two attack models within the ML supply chain (outsourcing training, retraining open-source model) to successfully backdoor different types of attacks (one-to-one attack and all-to-one attack). Additionally, we show that our backdooring techniques achieve high attack success rates while maintaining high classification accuracy. Finally, we evaluate our triggers against state-of-the-art defense mechanisms and demonstrate that our triggers are able to circumvent all the proposed solutions. Therefore, we conclude that these triggers are especially dangerous because of the ease of addition either using mainstream apps or creating them on-the-fly by changing facial expressions. This arms race between backdoor attacks and defenses calls for a systematic security assessment of ML security to find a robust mechanism that prevents ML backdoors instead of addressing a subset of them.\n\n\n\n\n\\section*{Appendix}\n\\subsection*{A. Attack detection for simple triggers}\n\\begin{figure}[h]\n \\centering\n \\includegraphics[width=0.4\\textwidth]{Imgs\/backdoor.png}\n \\caption{Top image: A genuine image from MNIST dataset. Bottom image: The same image with a dot trigger similar to Badnets \\cite{Badnets}. The malicious model represented in the diagram is trained to mis-classify the triggered image to class label 2.}\n \\label{fig:MNIST_images}\n\\end{figure}\n\n\\begin{figure*}\n \\centering\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.32\\textwidth]{Imgs\/MNIST_RegPCA_LR_new.png}\n \\label{fig:MNIST_pca_LR}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.32\\textwidth]{Imgs\/MNIST_L2Norm_LR_new.png}\n \\label{fig:MNIST_L2_LR}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.32\\textwidth]{Imgs\/MNIST_Eig_LR_new.png}\n \\label{fig:MNIST_Eig_Lr}\n }%\n \\caption{Backdoor analysis of MNIST triggers. \n \n (a) PCA of the datapoints considering the top two principle components at the learned representation level. The red dots are the malicious images of the targetted class and the green dots are the genuine images of the same class. (b) L2 norm of the genuine images (green) and malicious images (red) of the targetted label. (c) Correlation of the genuine (green) and malicious (red) images of the targetted label.} \n \\label{fig:MNIST_results}\n\\end{figure*}\n\n\\begin{figure*}\n \\centering\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Smile_activation.pdf}\n \\label{fig:smile_activation}\n \n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Arched_activation.pdf}\n \\label{fig:arched_activation}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Narrow_activation.pdf}\n \\label{fig:narrow_activation}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/Mouth_activation.pdf}\n \\label{fig:mouth_activation}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/YoungFaceApp_activation.pdf}\n \\label{fig:young_activation}\n \n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/OldFaceApp_activation.pdf}\n \\label{fig:old_activation}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/SmileFaceApp_activation.pdf}\n \\label{fig:smileFaceapp_activation}\n }%\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.12\\textwidth]{Imgs\/MakeupFaceApp_activation.pdf}\n \\label{fig:makeup_activation}\n }%\n \\caption{Activation analysis of different triggers for the neurons in the penultimate layer of the models. X-axis represents different neurons in the penultimate layers and y-axis represents the activation values of those neurons. Triggers: (a) Natural Smile (b) Arching of Eyebrows (c) Narrowing of eyes (d) Slightly opening of mouth (e) Young filter (f) Old filter (g) Smile filter (h) Makeup filter.} \n \\label{fig:activation_all}\n\\end{figure*}\n\nIn the Section \\ref{s:trigger_analysis} we evaluate our attacks to show that large, permeating and adaptive triggers are difficult to detect using state-of-the-art. In this section, we demonstrate the ease of detecting simple triggers for MNIST dataset. MNIST is a common dataset used by several backdoor attack and defense literature to evaluate their methodology \\cite{Badnets, NeuralCleanse,activation, Spectral_signatures}. We make an MNIST-BadNet using a one-pixel dot-trigger as shown in FIg. \\ref{fig:MNIST_images}.\n\nPrinciple Component Analysis (PCA) is transformation of data along the vectors of highest variance. These vectors, known as principle components, are orthogonal to each other and therefore, are linearly uncorrelated. Since the top components can define the data \\textit{sufficiently}, a very high dimensional data may be represented with its low-dimensional equivalent. However, we do not use PCA for dimensionality reduction\/ feature extraction. Rather we use it to find a reasonable representation of trigger samples and the genuine samples. For simple datasets like MNIST with distinct triggers, a simple PCA shows a distinction between malicious and genuine sub-populations. Different representations of images like normalized images, raw-data representation, or learned representation can be used to increase the effectiveness of PCA in distinguishing these sub-populations.\nWe use Euclidean distance (L2-norm) and correlations with top eigen vector as a measure to detect outliers (malicious sub-population) from the genuine data, as done for our triggers. In this sub-section we use the statistics used by literature to separate the sub-populations.\nL2 norm represents an image as a magnitude and images belonging to same class have similar L2 norms. Therefore, images that are slightly different albeit belonging to the same class, i.e. the malicious images, should have slightly different L2 norms. Fig. \\ref{fig:MNIST_L2_LR} shows that for MNIST with dot trigger, the sub-populations are easily separable by L2 norm at the learned representation level.\nTran et al. stated that correlation of images with the top Eigen vector of the dataset can be considered a spectral property of the malicious samples. This method of outlier detection is inspired from robust statistics. \nThe key intuition is that if the two sub-populations are distinguishable (using a particular image representation), then the malicious images along the direction of top Eigen vector will consist of a larger section of poisoned images. Therefore, they will have a different correlation than the genuine samples. \nFor simple datasets using small triggers (like described in Fig. \\ref{fig:MNIST_images}), the sub-populations are perfectly separable using correlation with top Eigen vector as shown in Fig. \\ref{fig:MNIST_Eig_Lr}.\n\\subsection*{B. Activation of different neurons}\nABS \\cite{CCS_ABS} detects backdoored models whose malicious behavior is successfully encoded with one neuron. In Fig. \\ref{fig:activation_all}, we show that the activation values peak for several neurons proving that our triggers are actually encoded using more than one neuron. Therefore, ABS would not be able detect our triggers.\n\n\\subsection*{C. Reversed triggers from Neural Cleanse}\nIn this section, we report the reversed engineered triggers for the backdoored models. Comparing the reversed triggers with Fig. \\ref{fig:social_media_filters}, we see stark differences between them. \n\n\\begin{figure*}\n \\centering\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.22\\textwidth]{Imgs\/FA_Smile_clean.pdf}\n \\label{fig:a}\n \n \\hspace{-6.5em}\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.22\\textwidth]{Imgs\/FA_Smile_fusion.pdf}\n \\label{fig:b}\n }%\n \\hspace{-6.5em}\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.22\\textwidth]{Imgs\/FA_Smile_re+cl.pdf}\n \\label{fig:c}\n }%\n \\hspace{-5em}\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.22\\textwidth]{Imgs\/Smile_clean.pdf}\n \\label{fig:d}\n }%\n \\hspace{-7.6em}\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.22\\textwidth]{Imgs\/Smile_fusion.pdf}\n \\label{fig:e}\n \n \\hspace{-7.6em}\n \\subfigure[]\n {%\n \n \\includegraphics[width=0.22\\textwidth]{Imgs\/Smile_re+cl.pdf}\n \\label{fig:f}\n }%\n \\caption{Neural cleanse defense analysis. Here we present the reversed triggers of 'smile' when embedded artificially using filters (a-c) and when embedded naturally using facial movements (d-f). (a) and (d) are the original images and (b) and (e) are the reverse-engineered triggers. In (c) and (f) we super impose the reversed triggers with the original images. As depicted, there is no visual similarity with the original triggers.} \n \\label{fig:NC_All}\n\\end{figure*}\n\n\n\n \n\n\\section{Introduction}\nIn the age of big data, the sheer volume of data discourages manual screening, and therefore, Machine learning has been an effective replacement to many conventional solutions. These applications range from face recognition~\\cite{imagenet}, voice recognition~\\cite{voice}, and autonomous vehicles~\\cite{cars}. Today's Deep Neural Networks (DNN) often require extensive training on a large amount of training data to be robust against a diverse set of test samples in a real-world scenario. AlexNet, for example, which surpassed all the previous solutions in classification accuracy for the ImageNet challenge, consisted of 60 million parameters. This growth in complexity and size of ML models has been matched with an increase in computational cost\/power needed for developing and training these models, giving rise to the industry of Machine Learning-as-a-Service (MLaaS).\n\nOutsourcing machine learning training meant an increase in the use of machine learning models by entities that would have otherwise been unable to develop and train complex models. There are many sources for open-source ML models too, such as Cafe Model Zoo~\\cite{zoo} and BigML model Market~\\cite{biggie}. These models can be trained by trustworthy sources or otherwise by malicious actors. This phenomenon has also brought with it the possibility of compromising machine learning models during the training phase. Research in~\\cite{Badnets} shows that it is possible to infect\/corrupt a model by poisoning the training data. This process introduces a backdoor\/trojan to the model. \nA backdoor in a DNN represents a set of neurons that are activated in the presence of unique triggers to cause malicious behavior.~\\cite{suppression} shows one such example where a dot (one-pixel trigger) can be used to trigger certain (backdoored\/infected) neurons in a model to maliciously change the true prediction of the model. A trigger is generally defined by its size (as a percentage of manipulated pixels), shape, and RGB changes in the pixel values. \nThere has been a plethora of proposed solutions that aim to defend against backdoored models through detection of backdoors, trigger reconstruction, and 'de-backdooring' infected models. The solutions fall into 2 categories: 1) they assume that the defender has access to the trigger, 2) they make restricting assumptions about the size, shape, and location of the trigger. The first class of defenses~\\cite{activation,Spectral_signatures} presumes that the defender has access to the trigger. Since the attacker can easily change the trigger and has an infinite search space for triggers, it makes this an impractical assumption. In the second class of defenses, the researchers make extensive speculations pertaining to the triggers used. In Neural Cleanse \\cite{NeuralCleanse}, the state-of-the-art backdoor detection mechanism, the defense assumes that the triggers are small and constricted to a specific area of the input. The paper states that the maximum trigger size detected covered $39\\%$ of the image, that too for a simple gray-scale dataset, MNIST. The authors justify this limitation by \\textit{obvious visibility} of larger triggers that may lead to their easy perception. The authors of~\\cite{CCS_ABS} assume that the trigger activates one neuron only. \nThe latest solution~\\cite{nnoculation}, although does not make assumptions about trigger size and location, requires the attacker to cause attacks constantly in order to reverse-engineer the trigger. \n\nFacial recognition is extensively used for biometric identification (in airports, ATMs, or offices) and cyber-crime\/forensics investigation field.\nThus, a possibility of trojanning such models may have catastrophic effects. Our triggers specifically target the facial recognition algorithms that are often the crux of \nautomatic detection of extremist posts~\\cite{abusive_content}, prevention of online dating frauds~\\cite{OnlineDating}, reporting of inappropriate images~\\cite{stealthy_images,Inapt} and airport immigration ~\\cite{immigration_uae}.\nWe study the impact of changes in the facial characteristics\/attributes in successfully trojanning a facial recognition model that was otherwise designed for an orthogonal task. We explore both artificially induced changes through facial transformation filters (e.g. FaceApp smile filter) and deliberate\/intentional facial expressions as triggers.\n\nIt is a matter of concern if simple filters can bypass all neural activations from the genuine features to maliciously drive the model to a mis-classification. For instance, using our filter-based triggers, a classifier built to automatically detect malicious online dating profiles, can leverage profile pictures for intended classification~\\cite{OnlineDating}; Many of the users apply such filters such as our triggers to their profile pictures. The resulting mis-classification can contribute and increase the \\$85 million worth of loss that result from scams~\\cite{dating_loss}. \nNatural triggers using facial muscle movement can be especially critical as many countries, such as the UAE, move towards ML based face recognition for automated passport control~\\cite{immigration_uae}. Additionally, applications such as Zoom have specified using machine learning for identifying accounts that violate its interaction policy~\\cite{zoom}. \nIn the attack literature, several types of patterns, like post-its on stop signs \\cite{Badnets}, specific black-rimmed spectacles \\cite{black_badnets} or specific patterns based on a desired masking size \\cite{NDSSTrojans} have been used to trigger backdoor neurons. Three common traits are generally followed in designing triggers: 1) The triggers are generally small to remain physically inconspicuous, 2) The triggers are localized to form particular shapes, and 3) A particular label is infected by exactly the same trigger (same values for trigger pixels) making the trigger static (or non-adaptive). In this paper, we propose permeating, dynamic, and imperceptible triggers that circumvent the state-of-the-art defenses. \n\n\n\n\nTo the best of our knowledge, this is the first attack to use specific facial gestures to trigger malicious behavior in facial recognition algorithms by constructing large scale, dynamic, and imperceptible triggers. We list our contributions as follows:\n\\begin{itemize}[nosep,leftmargin=1em,labelwidth=*,align=left]\n \\item We explore trojanning of ML models using filter-based triggers that artificially induce changes in facial attributes. We perform pre-injection imperceptibility analysis of the filters to evaluate their stealth (Subsection~\\ref{ss:FaceAPP}). \n \\item We study natural facial expressions as triggers to activate the maliciously-trained neurons to evaluate attack scenarios where trigger accessories (like glasses in case of Airport immigration) are not allowed (Subsection~\\ref{ss:natural}). \n \\item We explore different backdoor attack models on the ML supply chain. 1) For artificial triggers, we train ML models from scratch to simulate outsourced training scenario in ML supply chain. 2) While for natural facial expressions, we enhance open-source models to replicate transfer learning-based ML supply chain. Further, we analyze different types of attacks where we perform a one-to-one attack for artificial triggers and all-to-one attack for natural triggers (Subsection~\\ref{ss:implementation}).\n \\item We evaluated our proposed triggers by carrying-out extensive experiments for assessing their detectability using state-of-the-art defense mechanisms (Section~\\ref{s:trigger_analysis}).\n\\end{itemize}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nThe {\\sl International Gamma-ray Astrophysics Laboratory} ({\\sl INTEGRAL}) mission \\citep{2003A&A...411L...1W} has been observing the hard X-ray sky for more than 16 years. During this time it has uncovered a large number of new Be high-mass X-ray binaries (BeHMXBs), while also dramatically increasing the population of supergiant HMXBs (sgHMXBs; see e.g., Figure 4 in \\citealt{2015A&ARv..23....2W}). Additionally, {\\sl INTEGRAL} has unveiled new sub-classes of sgHMXBs: supergiant fast X-ray transients (SFXTs; \\citealt{2005A&A...444..221S,2006ESASP.604..165N}) and obscured supergiant HMXBs (e.g., IGR J16318-4848; \\citealt{2003MNRAS.341L..13M}). There are also several sub-classes of SFXTs, which exhibit different dynamical ranges in their X-ray luminosities, from $\\sim10^{2}-10^{5}$ and different flare time scales (see e.g., Table 2 in \\citealt{2015A&ARv..23....2W,2011ASPC..447...29C}). Among all of these systems there are also a number of ``intermediate'' HMXBs, that do not entirely fit into these categories and have properties that lie in between the classes (see e.g., \\citealt{2018MNRAS.481.2779S}). In order to understand the different physical processes responsible for producing the emission observed in these sub-classes and how they depend on the binary parameters (e.g., orbital period, eccentricity), it is important to increase the number of known members by identifying them through their X-ray and optical properties.\n\n\n\nPrior the to launch of {\\sl INTEGRAL}, BeHMXBs dominated the population of known HMXBs \\citep{2015A&ARv..23....2W}. BeHMXBs consist of a compact object (CO), typically a neutron star (NS), and a non-supergiant Be spectral type companion, generally having longer orbital periods, $P_{orb}\\gtrsim20$ days, than sgHMXBs (see \\citealt{2011Ap&SS.332....1R} for a review). The Be companions are rapidly rotating, have H$\\alpha$ emission lines, and have an infrared (IR) excess compared to B stars \\citep{2003PASP..115.1153P}. The H$\\alpha$ emission and IR excess come from a dense equatorial decretion disk rotating at or near Keplerian velocity, likely formed due to the fast rotation of the star (see \\citealt{2013A&ARv..21...69R} for a review). These systems can be both persistent, with X-ray luminosities on the order of $L_X\\approx10^{34}-10^{35}$ erg s$^{-1}$, or transient \\citep{2011Ap&SS.332....1R}. The transients become X-ray bright when the NS accretes material as it passes near (or through) the Be star's decretion disk. These systems exhibit two types of X-ray flares. Type I flares occur periodically or quasi-periodically when the NS passes through periastron, increasing the X-ray flux by an order of magnitude, and lasting a fraction of the orbital period ($\\sim0.2\\times P_{orb}\\approx1-10$ days; \\citealt{2011Ap&SS.332....1R}). Type II flares can occur at any orbital phase, reaching Eddington luminosity, and lasting a large fraction of an orbital period ($\\sim10-100$ days; \\citealt{2011Ap&SS.332....1R}). \n\nThe sgHMXBs are composed of a CO and a supergiant O or B spectral type companion, typically with shorter orbital periods, $P_{orb}\\lesssim10$ days, than BeHMXBs \\citep{2011Ap&SS.332....1R,2015A&ARv..23....2W}. If the CO orbits close enough to the companion star it can accrete via Roche lobe overflow, reaching X-ray luminosities up to $L_X\\sim10^{38}$ erg s$^{-1}$ during an accretion episode \\citep{2008A&A...484..783C,2015arXiv151007681C}. For longer period systems, the CO accretes from the fast ($\\sim1000$ km s$^{-1}$) radiative wind of the supergiant companion, leading to persistent X-ray luminosities of $L_X\\sim10^{35}-10^{36}$ erg s$^{-1}$ \\citep{2008A&A...484..783C}. These wind-fed systems are also often highly obscured ($N_H\\gtrsim10^{23}$ cm$^{-2}$) by the wind of the companion, and in some cases, by an envelope of gas and dust around the entire binary system (see e.g., IGR J16318-4848; \\citealt{2003A&A...411L.427W,2004ApJ...616..469F}). The supergiant stars in these systems also exhibit an H$\\alpha$ emission line due to their winds, which is often variable in shape and intensity and can have a P-Cygni profile (see e.g., Vela X-1, IGR J11215-5952; \\citealt{2001A&A...377..925B,2010ASPC..422..259L}).\n\n{\\sl INTEGRAL}'s wide field-of-view has enabled great progress in the study of HMXBs by discovering many SFXTs (see e.g., \\citealt{2006ESASP.604..165N}, or \\citealt{2013arXiv1301.7574S,2017mbhe.confE..52S} for reviews). Unlike the typical sgHMXBs, SFXTs exhibit much lower quiescent X-ray luminosities (as low as $L_X\\sim10^{32}$ erg s$^{-1}$) with highly energetic ($L_X\\sim10^{36}-10^{38}$ erg s$^{-1}$) X-ray flares lasting $\\sim100-10,000$ s (see e.g., \\citealt{2005A&A...441L...1I,2014A&A...568A..55R,2015A&A...576L...4R}). {\\sl INTEGRAL} has also uncovered several systems, with flare to quiescent X-ray luminosity ratios of $10^{2}-10^{3}$, lasting a few hours to days (see e.g., IGR J17354-3255; \\citealt{2011MNRAS.417..573S,2013A&A...556A..72D}). These systems have larger variability than seen in classic sgHMXB systems, but longer variability timescales than in SFXTs, and therefore have been called intermediate SFXTs \\citep{2011MNRAS.417..573S,2011ASPC..447...29C}.\n\n\nSeveral models have been put forward to describe the flaring phenomenon observed in SFXTs. One possibility is that the flares are caused by the accretion of an inhomogeneous clumpy wind, produced by the high-mass stellar companion, onto the compact object (see e.g., \\citealt{2005A&A...441L...1I,2006ESASP.604..165N,2007A&A...476..335W,2009MNRAS.398.2152D,2016A&A...589A.102B}). However, given the low-quiescent luminosities of some SFXTs, an additional mechanism for inhibiting the accretion onto the compact object is likely necessary (e.g., magnetic gating or sub-sonic accretion; \\citealt{2007AstL...33..149G,2008ApJ...683.1031B,2012MNRAS.420..216S}). \n\n\n\nThe source AX J1949.8+2534 was discovered by the {\\sl ASCA} Galactic plane survey having an absorbed flux of 6$\\times10^{-12}$ erg cm$^{-2}$ s$^{-1}$ in the 2-10 keV band \\citep{2001ApJS..134...77S}. AX J1949.8+2534 ( AX J1949 hereafter) was then detected by {\\sl INTEGRAL} for the first time in the hard X-ray band during two short flaring periods in 2015\/2016 \\citep{2015ATel.8250....1S}\\footnote{ The source is also referred to as IGR J19498+2534 in the {\\sl INTEGRAL} catalog of \\cite{2017MNRAS.470..512K}.}. The first flaring episode lasted $\\mathrel{\\hbox{\\rlap{\\hbox{\\lower4pt\\hbox{$\\sim$}}}\\hbox{$<$}}} 1.5$ days and reached a peak flux $F_X=1.1\\times10^{-10}$ erg cm$^{-2}$ s$^{-1}$ in the 22-60 keV band, while the second flaring episode lasted $\\sim4$ days with a similar peak flux $F_X= 1.0\\times10^{-10}$ erg cm$^{-2}$ s$^{-1}$ in the 22-60 keV band \\citep{2017MNRAS.469.3901S}. However, during the second flaring episode shorter time scale variability of $\\sim$2-8 ks was detected, reaching a peak flux of 2$\\times10^{-9}$ erg cm$^{-2}$ s$^{-1}$ in the 200 s binned light curve \\citep{2017MNRAS.469.3901S}. The dynamic range of the flare to quiescent X-ray luminosities in the 20-40 keV band is $\\gtrsim625$ \\citep{2017MNRAS.469.3901S}.\n\nIn soft X-rays, \\cite{2017MNRAS.469.3901S} also reported on a {\\sl Neil Gehrels Swift-XRT} observation of the source, which was detected with an absorbed flux of 1.8$\\times10^{-12}$ erg cm$^{-2}$ s$^{-1}$. This observation also provided a more accurate position for the source. However, there were two bright potential NIR counterparts within the 95\\% {\\sl Swift} positional uncertainty. Based on photometry, \\cite{2017MNRAS.469.3901S} classified these two bright NIR sources as B0V and B0.5Ia spectral type stars, respectively, leading to the conclusion that this source was either a BeHMXB or SFXT type source. The bright flares observed from AX J1949 disfavored the BeHXMB scenario, therefore, \\cite{2017MNRAS.469.3901S} favored the SFXT interpretation.\n\nIn this paper we report on {\\sl Neil Gehrels Swift-XRT}, {\\sl Chandra}, and {\\sl NuSTAR} legacy observations of the SFXT candidate AX J1949. These observations are part of an ongoing program to identify the nature of unidentified {\\sl INTEGRAL} sources by characterizing their broad-band (0.3-79.0 keV) X-ray spectrum. Additionally, we use the precise X-ray localization to identify the multi-wavelength counterpart for optical spectroscopic follow-up, which is also reported here. \n\n\n\n\\section{Observations and Data Reduction}\n\\label{obs_and_dat}\n\\subsection{{\\sl Neil Gehrels Swift X-ray Observatory}}\nThe {\\sl Neil Gehrels Swift} satellite's \\citep{2004ApJ...611.1005G} X-ray telescope (XRT; \\citealt{2005SSRv..120..165B}) has observed AX J1949 a total of six times, once in 2016 and five more times in a span of $\\sim20$ days in 2018. The 2016 observation was originally reported on by \\citealt{2017MNRAS.469.3901S} (although for consistency we reanalyze it here), while the five observations from 2018 are reported on for the first time in this paper. The details of each {\\sl Swift-XRT} observation can be found in Table \\ref{swiftobs}. We refer to each observation by the associated number shown in Table \\ref{swiftobs}. The data were reduced using HEASOFT version 6.25 and the 2018 July version of the {\\sl Swift-XRT} Calibration Database (CALDB). The {\\sl Swift-XRT} instrument was operated in photon counting mode and the event lists were made using the {\\tt xrtpipeline} task. To estimate the count rates, source significances, and upper limits we used the {\\tt uplimit} and source statistics tool ({\\tt sosta}) in the XIMAGE package (version 4.5.1). The count rates and upper limits were also corrected for vignetting, dead time, the point-spread function (PSF), and bad pixels\/columns that overlap with the source region. \n\nWe extracted the source spectrum for each {\\sl Swift} observation in which the source was detected with a signal-to-noise ratio (S\/N) $>3$ using XSELECT (version 2.4). A new ancillary response file was constructed using the {\\tt xrtmkarf} task using the exposure map created by the {\\tt xrtpipeline} task, and the corresponding response file was retrieved from the CALDB. The observed count rates (with no corrections applied) of the source in these observations were between $\\sim(8-12)\\times10^{-3}$ cts s$^{-1}$, so we used circular extraction regions with radii of 12 or 15 pixels depending on the observed count rate (see Table 1 in \\citealt{2009MNRAS.397.1177E}). Background spectra were extracted from a source free annulus centered on AX J1949. The {\\sl Swift} spectra were grouped to have at least one count per bin and fit using Cash statistics (C-stat; \\citealt{1979ApJ...228..939C}).\n\nIn this paper, all spectra are fit using XSPEC (version 12.10.1; \\citealt{1996ASPC..101...17A}). We used the Tuebingen-Boulder ISM absorption model ({\\tt tbabs}) with the solar abundances of \\cite{2000ApJ...542..914W}. All uncertainties reported in this paper (unless otherwise noted) are 1 $\\sigma$.\n\n\n\n\\begin{table}[h]\n\\caption{{\\sl Swift-XRT} Observations and signal-to-noise ratio of AX J1949.} \n\\label{swiftobs}\n\\begin{center}\n\\renewcommand{\\tabcolsep}{0.17cm}\n\\begin{tabular}{lcccccc}\n\\tableline \n$\\#$ & ObsID & offset &\tStart Time\t&\tExp. & S\/N \\\\\n\\tableline \n & & arcmin &\tMJD \t& s \\\\\n \\tableline \nS0$^{a}$ & 00034497001 & 2.422 & 57503.446 & 2932 & 5.3\\\\\nS1 &00010382001 & 2.997 & 58159.464 & 4910 & 5.5 \\\\\nS2 &00010382002 & 2.240 & 58166.446 & 3369 & 5.4 \\\\\nS3 &00010382003 & 2.441 & 58170.346& 1568 & 2.2\\\\\nS4 &00010382004 & 2.938 & 58173.336 & 4662 & 2.6 \\\\\nS5 &00010382005 & 0.424 & 58179.584 & 5482 & 4.5 \\\\\n\\tableline \n\\end{tabular} \n\\tablenotetext{a}{ This observation was originally reported on by \\cite{2017MNRAS.469.3901S}.}\n\\end{center}\n\\end{table}\n\n\n\n\\subsection{{\\sl Chandra} X-ray Observatory}\nWe observed AX J1949 using the {\\sl Chandra} Advanced CCD Imaging Spectrometer (ACIS; \\citealt{2003SPIE.4851...28G}) on 2018 February 25 (MJD 58174.547; obsID 20197) for 4.81 ks. The source was observed by the front-illuminated ACIS-S3 chip in timed exposure mode and the data were telemetered using the ``faint'' format. A 1\/4 sub-array was used to reduce the frame time to 0.8 s, ensuring that the pileup remained $<3\\%$ throughout the observation. All {\\sl Chandra} data analysis was performed using the {\\sl Chandra} Interactive Analysis of Observations (CIAO) software version 4.10 and CALDB version 4.8.1. The event file was reprocessed using the CIAO tool {\\tt chandra\\_repro} prior to analysis. \n\nTo locate all sources in the field of view, the CIAO tool {\\tt wavdetect} was run on the 0.5-8 keV band image\\footnote{ We ran {\\tt wavdetect} using the exposure map and PSF map produced by the CIAO tools {\\tt fluximage} and {\\tt mkpsfmap}, respectively.}. Only one source, the counterpart of AX J1949, was detected at the position R.A.$=297.48099^{\\circ}$, Dec.$=25.56639^{\\circ}$ with a statistical 95\\% confidence positional uncertainty of $0\\farcs13$ estimated using the empirical relationship (i.e., equation 12) from \\cite{2007ApJS..169..401K}. We were unable to correct for any systematic uncertainty in the absolute astrometry because only one X-ray source was detected. Therefore, we adopt the overall 90\\% {\\sl Chandra} systematic uncertainty of 0\\farcs8\\footnote{\\url{http:\/\/cxc.harvard.edu\/cal\/ASPECT\/celmon\/}} and convert it to the 95\\% uncertainty by multiplying by 2.0\/1.7. The statistical and systematic errors were then added in quadrature, giving a 95\\% confidence positional uncertainty radius of $0\\farcs95$ for the source.\n\nThe {\\sl Chandra} energy spectrum of AX J1949 was extracted from a circular region centered on the source and having a radius of 2$''$, enclosing $\\sim95\\%$ of the PSF at 1.5 keV\\footnote{See Chapter 4, Figure 4.6 at \\url{http:\/\/cxc.harvard.edu\/proposer\/POG\/html\/}.}, and containing 260 net counts. The background spectrum was extracted from a source free annulus centered on the source. Given the small number of counts, we fit the {\\sl Chandra} spectrum using Cash statistics (C-stat; \\citealt{1979ApJ...228..939C}).\n\n\n\nWe have also extracted the {\\sl Chandra} light curves using a number of different binnings to search for spin and orbital periods in the data. Prior to extraction, the event times were corrected to the Solar system barycenter using the CIAO tool {\\tt axbary}.\n\n\n\\subsection{{\\sl NuSTAR}}\nThe {\\sl Nuclear Spectroscopic Telescope Array} ({\\sl NuSTAR}; \\citealt{2013ApJ...770..103H}) observed AX J1949 on 2018 February, 24 (MJD 58173.070; obsID 30401002002) for 45 ks. We reduced the data using the NuSTAR Data Analysis Software (NuSTARDAS) version 1.8.0 with CALDB version 20181022. Additionally, we filtered the data for background flares caused by {\\sl NuSTAR's} passage through the South Atlantic Anomaly (SAA) using the options {\\tt saacalc}=2, {\\tt saamode}=optimized, and {\\tt tentacle}=yes, which reduced the total exposure time to 43 ks. \n\n\nThe source's energy spectra from the FPMA and FPMB detectors were extracted from circular regions with radii of 45$''$ centered on the source position. The background spectra were extracted from source-free regions away from AX J1949, but on the same detector chip. The {\\sl NuSTAR} spectra were grouped to have at least one count per bin. {\\sl NuSTAR} light curves were also extracted from both the FPMA and FPMB detectors in the 3-20 keV energy range using a number of different bin sizes (i.e., 100 s, 500 s, 1 ks, 5 ks). All {\\sl NuSTAR} light curves plotted in this paper show the averaged (over the FPMA and FPMB detectors) net count rate.\n\n\\subsection{MDM Spectroscopy}\n\nThe accurate position of the X-ray counterpart to AX J1949 provided by {\\sl Chandra} has allowed us to identify the optical\/NIR counterpart to the source. The optical\/NIR counterpart is the brightest source (see Table \\ref{mw_mag}) considered by \\cite{2017MNRAS.469.3901S} as a potential counterpart (i.e., their source 5). This source has a {\\sl Gaia} position R.A.$=297.480949745(9)^{\\circ}$ and Decl.$=25.56659555(1)^{\\circ}$, which is $\\sim0\\farcs76$ away from the {\\sl Chandra} source position and within the 2$\\sigma$ {\\sl Chandra} positional uncertainty.\n\nOn 2018 October 18, a 600~s spectrum was obtained with the Ohio State\nMulti-Object Spectrograph on the 2.4~m Hiltner telescope of the\nMDM Observatory on Kitt Peak, Arizona. A $1.\\!^{\\prime\\prime}2$\nwide slit and a volume-phase holographic grism provided a dispersion\nof 0.72 \\AA\\ pixel$^{-1}$ and a resolution of $\\approx3$ \\AA\\ over\nthe wavelength range 3965--6878 \\AA. The reduced spectrum is\nshown in Figure \\ref{opt_spec}, where it can be seen that flux is not well\ndetected below 4900 \\AA\\ due to the large extinction to the star.\nAlthough a standard star was used for flux calibration, the narrow\nslit and partly cloudy conditions are not conducive to\nabsolute spectrophotometry. Therefore, as a last step we have \nscaled the flux to match the $V$ magnitude from Table \\ref{mw_mag}.\n\n\\begin{figure*}\n\\centering\n\\includegraphics[trim={0 0 0 0},scale=0.60]{Figure1.pdf}\n\\caption{MDM 2.4 m optical spectrum of AX J1949 showing H$\\alpha$ emission and He I absorption. The other spectral absorption features are diffuse interstellar bands caused by the large reddening towards the source ($A_V=$8.5-9.5, see Section \\ref{spec_type}). (Supplemental data for this figure are available in the online journal.)\n\\label{opt_spec}\n}\n\\end{figure*}\n\n\n\\section{Results}\n\n\n\\subsection{X-ray spectroscopy}\n\\label{xspectro}\n AX J1949 is a variable X-ray source (see Figure \\ref{xrt_cr}), but fortunately, {\\sl Chandra} observed it during a relatively bright state, allowing for the most constraining X-ray spectrum of all of the observations reported here. We fit the {\\sl Chandra} spectrum in the 0.5-8 keV range with two models, an absorbed power-law model and an absorbed blackbody model (see Figure \\ref{cxo_spec}). The best-fit power law model has a hydrogen absorption column density $N_{H}=7.5^{+1.6}_{-1.4}\\times10^{22}$ cm$^{-2}$, photon index $\\Gamma=1.4\\pm0.4$, and absorbed flux $(2.0\\pm0.3)\\times10^{-12}$ erg cm$^{-2}$ s$^{-1}$ (in the 0.3-10 keV band) with a C-stat of 121.5 for 182 d.o.f. On the other hand, the best-fit blackbody model has hydrogen absorption column density $N_{H}=(4.7\\pm1)\\times10^{22}$ cm$^{-2}$, temperature $kT=1.5\\pm0.2$ keV, radius $r_{\\rm BB}=150^{+30}_{-20}d_{\\rm 7kpc}$ m, and absorbed flux $(1.6\\pm0.2)\\times10^{-12}$ erg cm$^{-2}$ s$^{-1}$ (in the 0.3-10 keV band), where $d_{\\rm 7kpc}=d\/$(7 kpc) is the assumed distance to the source (see Section \\ref{spec_type}), with a C-stat of 119.3 for 182 d.o.f. In either case, the photon-index, or the blackbody temperature and emitting radius, are consistent with those observed in SFXTs (see e.g., Table 4 in \\citealt{2009MNRAS.399.2021R}). Therefore, since both models fit the data about equally well we cannot distinguish between them.\n\n\\begin{figure}\n\\centering\n\\includegraphics[trim={0 0 0 0},scale=0.38]{Figure2.pdf}\n\\caption{{\\sl Swift-XRT} net count rate light curve of AX J1949 during the $\\sim$20 day observing period in 2018, which included observations S1 through S5 (see Table \\ref{swiftobs}). \n\\label{xrt_cr}\n}\n\\end{figure}\n\n\n\n\n\n\nWe extracted the {\\sl Swift-XRT} spectra for all observations where the source was detected with a S\/N$>3$ (i.e., S0, S1, S2, and S5). The source is not bright enough in any of these single observations to fit a constraining spectrum. Therefore, we jointly fit the {\\sl Chandra} spectrum with the four {\\sl Swift-XRT} spectra, tying together the hydrogen absorbing column, but leaving the photon index and normalization free to vary for each spectrum. The best-fit parameters can be seen in Table \\ref{spec_par}. There is marginal evidence ($\\sim2\\sigma$) of spectral softening in the spectrum of S2, which could also be due to a variable hydrogen absorbing column as is often seen in SFXTs (see e.g., \\citealt{2009MNRAS.397.1528S,2016A&A...596A..16B,2018arXiv181111882G}). We have also fit the spectra with an absorbed power-law model after tying together the $N_H$ and photon-index before fitting. When this is done, the best fitting spectrum has a hydrogen absorption column density $N_H=5.6^{+1.1}_{-1.0}\\times10^{22}$ cm$^{-2}$ and photon-index $\\Gamma=1.2\\pm0.3$ with a C-stat of 275.0 for 295 d.o.f.\n\n\\begin{figure}\n\\centering\n\\includegraphics[trim={0 0 0 0},scale=0.41]{Figure3.pdf}\n\\caption{{\\sl Chandra} 0.5-8 keV spectrum with best-fit power-law model (see Section \\ref{xspectro} for best-fit parameters). The middle panel shows the ratio of the data to the power-law model, while the bottom panel shows the ratio of the data to the blackbody model. The spectra were binned for visualization purposes only.\n\\label{cxo_spec}\n}\n\\end{figure}\n\nUnfortunately, AX J1949 was simultaneously observed with {\\sl NuSTAR} during the S4 {\\sl Swift-XRT} observation, when the source was at its faintest flux level (see Figure \\ref{fluxes}). Although the source was still jointly detected (i.e., FPMA+FPMB) with a significance of 5.3$\\sigma$ and 5.7$\\sigma$ (in the 3-79 keV and 3-20 keV energy bands, respectively), there were too few counts to constrain its spectrum. However, we extracted the 3$\\sigma$ upper-limit flux in the 22-60 keV energy range for later comparison to {\\sl INTEGRAL}. To do this, we calculated the 3$\\sigma$ upper-limit on the net count rate of the source in the 22-60 keV energy range (0.0011 cts s$^{-1}$). SFXTs often show a cutoff in their spectra between $\\sim5-20$ keV (see e.g., \\citealt{2011MNRAS.412L..30R,2017ApJ...838..133S}) so we use a power-law model with the best-fit absorption ($N_H=5.6\\times10^{22}$ cm$^{-2}$) and a softer photon index ($\\Gamma=2.5$) that is typical for SFXTs detected in hard X-rays by {\\sl INTEGRAL} (see e.g., \\citealt{2008A&A...487..619S,2011MNRAS.417..573S}) to convert the count rate to flux. The 3$\\sigma$ upper-limit on the flux in the 22-60 keV energy range is $F=5.8\\times10^{-13}$ erg cm$^{-2}$ s$^{-1}$.\n\n\n\n\\begin{table}[t!]\n\\caption{Photon indices for the simultaneous best-fit absorbed power-law model.} \n\\label{spec_par}\n\\begin{center}\n\\renewcommand{\\tabcolsep}{0.11cm}\n\\begin{tabular}{lccccc}\n\\tableline \nObs. & $N_H$ & $\\Gamma$ &\t\tC-stat\t &d.o.f. \\\\\n\\tableline \n & $10^{22}$ cm$^{-2}$ &\t \t & \\\\\n\\tableline \nCXO & 5.6$^{+1.1}_{-1.0}$ & 1.0$\\pm0.3$ & 243.4 & 286 \\\\ \nS0 & 5.6\\tablenotemark{a} & 1.0$\\pm0.5$ & -- & -- \\\\\nS1 & 5.6\\tablenotemark{a} & 1.5$\\pm0.6$ & -- & -- \\\\\nS2 & 5.6\\tablenotemark{a} & 2.5$^{+0.6}_{-0.5}$ & -- & --\\\\\nS5 & 5.6\\tablenotemark{a} & 2.2$\\pm0.7$ & -- & --\\\\\n\\tableline \n\\end{tabular} \n\\tablenotetext{a}{The hydrogen absorption column density ($N_{\\rm H}$) values were tied together for the fit reported in this table.}\n\\end{center}\n\\end{table}\n\n\n\n\\subsection{X-ray variability and timing}\nThe fluxes of the source are not constrained if the photon-index is left free when jointly fitting the {\\sl Chandra} and {\\sl Swift-XRT} spectra. Therefore, to extract the fluxes from the {\\sl Swift} and {\\sl Chandra} observations we use the jointly fit model (i.e., $\\Gamma=1.2$ and $N_H=5.6\\times10^{22}$ cm$^{-2}$; see Section \\ref{xspectro}). The fluxes derived from these fits can be seen in Figure \\ref{fluxes}. Additionally, we assumed the same best-fit power-law model when converting the {\\sl Swift} 3$\\sigma$ upper-limit count rates to fluxes. The flux from observation S0 is not shown in this figure but is $F_{0.3-10 \\ keV}=(1.4\\pm0.3)\\times10^{-12}$ erg cm$^{-2}$ s$^{-1}$, which is consistent within 2$\\sigma$ of the value reported by \\cite{2017MNRAS.469.3901S}.\n\n\\begin{figure}\n\\centering\n\\includegraphics[trim={0 0 0 0},scale=0.36]{Figure4.pdf}\n\\caption{Fluxes of AX J1949 in the 0.3-10 keV energy range from the jointly fit {\\sl Swift-XRT} and {\\sl Chandra} (CXO) spectra (see Section \\ref{xspectro}). The points show the fluxes from the different observations (circles for {\\sl Swift-XRT} and squares for {\\sl Chandra}), while the triangles show the 3$\\sigma$ upper-limit of the source's flux when it was not detected in the {\\sl Swift-XRT} observations. The flux for observation S0 is excluded. The {\\sl NuSTAR} observation was concurrent with the S4 observation (i.e., the second triangle), when the flux of the source was at its minimum.\n\\label{fluxes}\n}\n\\end{figure}\n\n\nWe have also searched for shorter timescale variability of AX J1949 in the individual {\\sl Chandra} and {\\sl NuSTAR} observations. The 1 ks binned net count rate light curve from the {\\sl Chandra} observation is shown in Figure \\ref{cxo_cr}. The light curve clearly shows that the source brightness decreased throughout the observation on a several ks timescale. No significant flaring behavior was found in the {\\sl Chandra} light curves with smaller binnings. We also searched for any periodic signal in the {\\sl Chandra} data using the Z$^{2}_{1}$ \\citep{1983A&A...128..245B}, but no significant periodicity was found.\n\n\\begin{figure}\n\\centering\n\\includegraphics[trim={0 0 0 0},scale=0.3]{Figure5.pdf}\n\\caption{{\\sl Chandra} 0.5-8 keV net count rate during the 4.81 ks observation with a 1 ks binning. The source was variable on $\\sim1$ ks timescales during this observation.\n\\label{cxo_cr}\n}\n\\end{figure}\n\nThere are also indications of variability in the 3-20 keV {\\sl NuSTAR} light curve. Notably, the 5 ks binned net light curve (see Figure \\ref{nus_cr}) shows evidence of a flux enhancement lasting $\\sim$10 ks. To test the significance of this variability we have fit a constant value to the light curve. The best-fit constant has a value of 3.0$\\pm0.9\\times10^{-3}$ counts s$^{-1}$ with a $\\chi^2$=93.6 for 17 degrees of freedom, implying that the source is variable at a $\\gtrsim6.5\\sigma$ level. The 500 s binned light curve of the period of flux enhancement (see Figure \\ref{ns_flare}) shows the light curve reached a peak count rate of $\\sim$0.03 cts s$^{-1}$. This corresponds to a peak flux of $\\sim1.5\\times10^{-12}$ erg cm$^{-2}$ s$^{-1}$, assuming the best-fit power law model (i.e., $\\Gamma=1.2$ and $N_{H}=5.6\\times10^{22}$ cm$^{-2}$). This flux is consistent with the flux observed during the {\\sl Chandra} observation, suggesting that the source's light curve varies over timescales of $\\sim1-10$ ks.\n\n\n\n\n\n\\begin{figure}\n\\centering\n\\includegraphics[trim={0 0 0 0},scale=0.37]{Figure6.pdf}\n\\caption{{\\sl NuSTAR} 3-20 keV averaged (over the FPMA and FPMB detectors) net count rate light curve with a 5 ks binning. The source was variable on $\\sim5$ ks timescales during this observation, including a period of flux enhancement lasting $\\sim10$ ks. The gray band shows the time and duration of the joint {\\sl Swift-XRT} observation (observation S4). \n\\label{nus_cr}\n}\n\\end{figure}\n\n\n\n\\begin{figure}\n\\centering\n\\includegraphics[trim={0 0 0 0},scale=0.37]{Figure7.pdf}\n\\caption{{\\sl NuSTAR} 3-20 keV net count rate light curve with a 500 s binning during the period of flux enhancement. The source reaches a peak net count-rate of $\\sim$0.03 cts s$^{-1}$. \n\\label{ns_flare}\n}\n\\end{figure}\n\n\n\n\n\\subsection{Optical\/NIR companion photometry}\n\nUsing the accurate source position provided by {\\sl Chandra} we have confidently identified the multi-wavelength counterpart to this source. Only one {\\sl Gaia} source lies within the $\\sim1''$ (2$\\sigma$) positional uncertainty of AX J1949. The source is also detected in optical by Pan-STARRs \\citep{2016arXiv161205243F}, in NIR by 2MASS \\citep{2006AJ....131.1163S}, and in IR by {\\sl WISE} \\citep{2012wise.rept....1C}. The counterparts multi-wavelength magnitudes can be found in Table \\ref{mw_mag}.\n\n\n\\begin{table}[t!]\n\\caption{Multi-wavelength magnitudes of the counterpart to AX J19498.} \n\\label{mw_mag}\n\\small\n\\begin{center}\n\\renewcommand{\\tabcolsep}{0.07cm}\n\\begin{tabular}{lcccc}\n\\tableline \nObs. & R.A. & Decl. & Filter & Mag. \\\\\n\\tableline \n & deg. & deg. & &\\\\\n\\tableline\n{\\sl Gaia} & 297.480949745(9) & 25.56659555(1) & G & 13.9860(8) \\\\\n & -- & -- & Bp & 16.299(6) \\\\\n & -- & -- & Rp & 12.578(3) \\\\\nPS\\tablenotemark{a} & 297.480960(2) & 25.566612(7) & g & 17.516(8) \\\\ \n & -- & -- & r & 14.981(4) \\\\\n & -- & -- & i & 13.1990 \\\\\n & -- & -- & z & 12.1140 \\\\\n & -- & -- & y & 11.7220 \\\\\n NOMAD & 297.4809331 & 25.5666150 & B & 18.04 \\\\\n & -- & -- & V & 16.12 \\\\\n & -- & -- & R & 14.62 \\\\\n 2MASS & 297.480989 & 25.566639 & J & 9.90(2) \\\\\n & -- & -- & H & 9.07(2) \\\\\n & -- & -- & $K_s$ & 8.63(2) \\\\\n {\\sl WISE} & 297.4809572 & 25.5666135 & W1 & 8.32(2) \\\\\n & -- & -- & $W2$ & 8.17(2) \\\\\n & -- & -- & $W3$ & 8.30(3) \\\\\n \\tableline \n\\end{tabular} \n\\tablenotetext{a}{Pan-Starrs}\n\\end{center}\n\\end{table}\n\n{\\sl Gaia} has also measured the parallax, $\\pi=0.081(64)$, and proper motion, $\\mu_{\\alpha}\\cos{\\delta}=-2.37(9)$, $\\mu_{\\delta}=-5.2(1)$, of AX J1949's counterpart \\citep{2018A&A...616A...1G}. Unfortunately, the {\\sl Gaia} parallax has a large relative error ($\\sim80\\%$) and, therefore, has a large uncertainty in the inferred distance, $d=6.1^{+2.3}_{-1.5}$ kpc \\citep{2018AJ....156...58B}. We note that this source also has a poorly fit astrometric solution with excess noise, possibly due to its binary nature. Therefore, we do not rely on the {\\sl Gaia} parallax and simply take it as an indicator that the source is at least at a distance of a few kpc. The large extinction observed in the optical magnitudes of the source also support this claim (see Section \\ref{spec_type}). \n\n\n\\subsection{Optical spectroscopy}\n\nThe optical counterpart of AX J1949 shows H$\\alpha$ emission with \nEW$\\approx -1.6$ \\AA\\ and FWZI $\\approx500$ km~s$^{-1}$, and absorption\nlines of \\ion{He}{1} 5876 \\AA\\ and 6678 \\AA\\ (see Figure \\ref{opt_spec}). The remaining absorption\nfeatures are diffuse interstellar bands, whose strengths\nare consistent with the large extinction inferred from the colors\nof the star (see Section \\ref{spec_type}).\n\n\n\n\n\n\n\\section{Discussion}\n\n\\subsection{Optical\/NIR Companion Stellar Type}\n\\label{spec_type}\nThe optical spectrum of the companion star shows He I absorption lines, suggesting that the star is hot (i.e., $\\mathrel{\\hbox{\\rlap{\\hbox{\\lower4pt\\hbox{$\\sim$}}}\\hbox{$>$}}}$10,000 K) and, therefore, is likely to be of either O or B spectral class. Here we use the NOMAD $V$-band photometry in combination with the 2MASS NIR photometry to assess the spectral type, extinction, and distance to this star. The NOMAD $V$-band magnitude and the 2MASS NIR magnitudes can be found in Table \\ref{mw_mag}. To estimate the extinction we use the intrinsic color relationships between the $V$-band and NIR magnitudes provided by Table 2 in \\cite{2014AcA....64..261W}. To recover the intrinsic color relationships of late O and early to late B type stars, we need to deredden the source by an $A_V\\approx8.5-9.5$, assuming the extinction law of \\cite{1989ApJ...345..245C}. The 3D extinction maps\\footnote{We use the {\\tt mwdust} python package to examine the dust maps \\citep{2016ApJ...818..130B}.} of \\cite{2006A&A...453..635M} suggest a distance of $\\approx$ 7-8 kpc for this range of $A_V$.\n\n\n\n\nThe reddening and estimated distance to the source can be used to place constraints on its spectral type and luminosity class. In order for the star to have an observed $V$-band magnitude of 16.12 at a distance of 7 kpc with a reddening of $A_V=8.5$, it should have an absolute magnitude of $M_V\\approx-6.6$. The lack of He II absorption features in the optical spectrum suggests that the star is more likely to be of B-type rather than O-type. By comparing the absolute magnitude to Table 7 in \\cite{2006MNRAS.371..185W} we find that the star is most likely to be of the Ia luminosity class. Further, the strong He I absorption lines indicate that the star is an early B-type star. However, we mention one caveat, which is that while the lower luminosity classes are not bright enough to be consistent with the distance\/reddening to the source, there is a large uncertainty in some of the absolute magnitudes for different luminosity classes (e.g., luminosity class Ib, II; \\citealt{2006MNRAS.371..185W}). \n\n\n\nLastly, the optical spectrum of the star shows an H$\\alpha$ emission line, which is seen in both Be, and supergiant type stars. This line has been observed in a number of similar systems and is often variable in shape and intensity, sometimes showing a P-Cygni like profile (see e.g., \\citealt{2006ApJ...638..982N,2006ESASP.604..165N,2006A&A...455..653P}). Unfortunately, we have only obtained a single spectrum and cannot assess the variability of the H$\\alpha$ line. Additionally, it does not appear to be P-Cygni like in this single spectrum. \n\n\n\n\nIn comparison, by using optical and NIR photometry, \\cite{2017MNRAS.469.3901S} found the stellar companion of AX J1949 to be consistent with a B0.5Ia type star at a distance, $d=8.8$ kpc, and having a reddening, $A_V=7.2$. We have also found that the star is consistent with an early B-type type star of the Ia luminosity class. However, our estimates of the reddening and distance differ slightly from \\cite{2017MNRAS.469.3901S}, in that we find a larger reddening and smaller distance. \n\n\n\\subsection{SFXT Nature}\n\nAt an assumed distance of $\\sim7$ kpc, AX J1949's {\\sl Chandra} and {\\sl Swift-}XRT fluxes (when the source is significantly detected) correspond to luminosities of $\\sim1\\times10^{34}$ erg s$^{-1}$ and $\\sim3-8\\times10^{33}$ erg s$^{-1}$, respectively. Further, at harder X-ray energies, AX J1949 has a {\\sl NuSTAR} 3$\\sigma$ upper-limit quiescent luminosity $\\mathrel{\\hbox{\\rlap{\\hbox{\\lower4pt\\hbox{$\\sim$}}}\\hbox{$<$}}} 3\\times10^{33}$ erg s$^{-1}$ in the 22-60 keV band, while during its flaring X-ray activity detected by {\\sl INTEGRAL}, it reached a peak luminosity of $\\sim10^{37}$ erg s$^{-1}$ in the 22-60 keV band \\citep{2017MNRAS.469.3901S}. This implies a lower-limit on the dynamical range of the system of $\\mathrel{\\hbox{\\rlap{\\hbox{\\lower4pt\\hbox{$\\sim$}}}\\hbox{$>$}}} 3000$. The large dynamical range and relatively low persistent X-ray luminosity are typical of SFXTs. Additionally, converting the absorbing column density, $N_{\\rm H}=5.6\\times10^{22}$ cm$^{-2}$, derived from the spectral fits to the X-ray data to $A_V$ following the relationship provided by \\cite{2015MNRAS.452.3475B}, we find $A_V\\approx20$. This $A_V$ is a factor of 2-2.5 larger than the $A_V$ derived from the interstellar absorption to the companion star (see Section \\ref{spec_type}). This suggests that a large fraction of the absorption is intrinsic to the source, which has also been observed in other SFXTs (e.g., XTE J1739-302, AX J1845.0-0433; \\citealt{2006ApJ...638..982N,2009A&A...494.1013Z}). Lastly, the optical counterpart to AX J1949 appears to be an early supergiant B-type star, which have been found in $\\sim40\\%$ of the SFXT systems (see Figure 2 in \\citealt{2017mbhe.confE..52S}). NIR spectra and additional optical spectra should be undertaken to search for additional absorption features, as well as changes in the H$\\alpha$ line profile to solidify the spectral type. \n\n\n\n\n\n\n\n\\section{Summary and Conclusion}\n\nA large number of unidentified {\\sl INTEGRAL} sources showing rapid hard X-ray variability have been classified as SFXTs. We have analyzed {\\sl Neil Gehrels Swift-XRT}, {\\sl Chandra}, and {\\sl NuSTAR} observations of the SFXT candidate AX J1949, which shows variability on ks timescales. The superb angular resolution of {\\sl Chandra} has allowed us to confirm the optical counterpart to AX J1949 and obtain its optical spectrum, which showed an H$\\alpha$ emission line, along with He I absorption features. The spectrum, coupled with multi-wavelength photometry allowed us to place constraints on the reddening, $A_V=8.5-9.5$, and distance, $d\\approx7-8$ kpc, to the source. We find that an early B-type Ia is the most likely spectral type and luminosity class of the star, making AX J1949 a new confirmed member of the SFXT class.\n\n\n\\medskip\\noindent{\\bf Acknowledgments:}\nWe thank Justin Rupert for obtaining the optical spectrum at MDM. We thank the anonymous referee for providing useful and constructive comments that helped to improve the paper. This {\\sl work} made use of observations obtained at the MDM Observatory, operated by Dartmouth College, Columbia University, Ohio State University, Ohio University, and the University of Michigan. This work made use of data from the {\\it NuSTAR} mission, a project led by the California Institute of Technology, managed by the Jet Propulsion Laboratory, and funded by the National Aeronautics and Space Administration. We thank the {\\it NuSTAR} Operations, Software and Calibration teams for support with the execution and analysis of these observations. This research has made use of the {\\it NuSTAR} Data Analysis Software (NuSTARDAS) jointly developed by the ASI Science Data Center (ASDC, Italy) and the California Institute of Technology (USA). JH and JT acknowledge partial support from NASA through Caltech subcontract CIT-44a-1085101. JAT acknowledges partial support from Chandra grant GO8-19030X.\n\n \\software{CIAO (v4.10; \\citealt{2006SPIE.6270E..1VF}), XSPEC (v12.10.1; \\citealt{1996ASPC..101...17A}), NuSTARDAS (v1.8.0), Matplotlib \\citep{2007CSE.....9...90H}, Xselect (v2.4e), XIMAGE (v4.5.1), HEASOFT (v6.25), MWDust \\citep{2016ApJ...818..130B}}\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\n The study of fractional order differential equations has been attracting many scientists because of its adequate and interesting applications in modeling of real-life problems related to several fields of science \\cite{L1}-\\cite{L5}. Initial-value problems (IVPs) and boundary-value problems (BVPs) involving the Riemann-Liouville and Caputo derivatives attract most interest (see, for instance, \\cite{L6}, \\cite{L7}, \\cite{L8}). Especially, studying IVPs and BVPs for the sub-diffusion, fractional wave equations are well-studied (see \\cite{sad}, \\cite{ruzh}, \\cite{karruzh}). BVPs for mixed type equations are also an interesting target for many authors (see \\cite{L9}-\\cite{L13}).\n\n Introducing a generalized Riemann-Liouville fractional derivatives (it is called Hilfer's deriva\\-tive) has opened a new gate in the research of fractional calculus (\\cite{L14}-\\cite{L16}). Therefore, one can find several works devoted to studying this operator in various problems \\cite{L17}, \\cite{L18}. We also note that in 1968, M.~M.~Dzhrbashyan and A.~B.~Nersesyan introduced the following integral-differential operator \\cite{Ldn}\n\\begin{equation}\\label{DN}\nD_{0x}^{\\sigma_n}g(x)=I_{0x}^{1-\\gamma_n}D_{0x}^{\\gamma_{n-1}}...D_{0x}^{\\gamma_{1}}D_{0x}^{\\gamma_{0}}g(x), ~ ~ n\\in\\mathbb{N}, ~ x>0,\n\\end{equation}\nwhich is more general than Hilfer's operator. Here $I_{0x}^{\\alpha}$ and $D_{0x}^{\\alpha}$ are the Riemann-Liouville fractional\nintegral and the Riemann-Liouville fractional derivative of order $\\alpha$ respectively (see Definition 2.1), $\\sigma_n\\in(0, n]$ which is defined by\n$$\n\\sigma_n=\\sum\\limits_{j=0}^{n}\\gamma_j-1>0, ~ \\gamma_j\\in(0, 1].\n$$\nThere are some works \\cite{bag}, \\cite{bag2}, related with this operator. New wave of researches involving this operator might appear due to the translation of original work \\cite{Ldn} in FCAA \\cite{Lfca}.\n \n In addition, from announcing the concept of hyper-Bessel fractional differential derivative by I. Dimovski \\cite{L21}, several articles have been published dedicated to studying problems containing this type of operators (see \\cite{L22}-\\cite{L26}). For instance, fractional diffusion equation and wave equation were widely investigated in different domains in \\cite{L19}-\\cite{L20}.\n\nIn this work, we investigate a boundary value problem for a mixed equation involving the sub-diffusion equation with Caputo-like counterpart of a hyper-Bessel fractional differential operator and the fractional wave equation with Hilfer's bi-ordinal derivative in a rectangular domain. The theorem about the uniqueness and existence of the solution is proved.\n\n{The rest of the paper is organized as follows: In Preliminaries section we provide necessary information on Mittag-Leffler functions (Section 2.1), hyper-Bessel functions (Section 2.2.), bi-ordinal Hilfer's fractional derivatives (Section 2.3) and on differential equation involving bi-ordinal Hilfer's fractional derivatives (Section 2.4). Auxiliary result is formulated in Theorem 2.2. In Section 3, we formulate the main problem and state our main result in Theorem 3.1. In Appendix one can find detailed arguments of the proof of Theorem 2.1.}\n\n\n\\section{Preliminaries}\n\nIn this section we present some definitions and auxiliary results related to generalized Hilfer's derivative and fractional hyper-Bessel differential operator which will be used in the sequel. We start recalling the definition of the Mittag-Leffler function.\n\n\\subsection{Important properties of the Mittag-Leffler function}\n\nThe two parameter Mittag-Leffler (M-L) function is an entire function given by\n\\begin{equation}\\label{e1}\nE_{\\alpha, \\beta}(z)=\\sum_{k=0}^{\\infty}\\frac{z^k}{\\Gamma(\\alpha{k}+\\beta)}, \\, \\, \\, \\alpha>0, \\, \\beta\\in\\mathbb{R}.\n\\end{equation}\n\n\\textbf{Lemma 2.1} (see \\cite{L3}) Let $\\alpha<2, \\, \\beta\\in\\mathbb{R}$ and $\\frac{\\pi\\alpha}{2}<\\mu\\alpha$ and $x\\geq0$ one has\n$$\n\\frac{1}{\\Big(1+\\sqrt{\\frac{\\Gamma(1-\\alpha)}{\\Gamma(1+\\alpha)}}x\\Big)^2}\\leq{E_{\\alpha, \\alpha}(-x)}\\leq\\frac{1}{\\Big(1+\\frac{\\Gamma(1+\\alpha)}{\\Gamma(1+2\\alpha)}x\\Big)^2}\n$$\nand\n$$\n\\frac{1}{1+\\frac{\\Gamma(\\beta-\\alpha)}{\\Gamma(\\beta)}x}\\leq\\Gamma(\\beta){E_{\\alpha, \\beta}(-x)}\\leq\\frac{1}{1+\\frac{\\Gamma(\\beta)}{\\Gamma(\\beta+\\alpha)}x}.\n$$\n\nThe Laplace transform of M-L function is given in the following lemma.\n\n\\textbf{Lemma 2.2.} (\\cite{L8}) For any $\\alpha>0, \\, \\, \\beta>0$ and $\\lambda\\in\\mathbb{C}$, we have\n$$\n\\mathcal{L}\\{t^{\\beta-1}E_{\\alpha, \\beta}(\\lambda{t}^\\alpha)\\}=\\frac{s^{\\alpha-\\beta}}{s^\\alpha-\\lambda}, \\, \\, (Re(s)>\\mid{\\lambda}\\mid^{1\/\\alpha}),\n$$\nwhere the Laplace transform of a function $f(t)$ is defined by\n\n$$\n\\mathcal{L}\\{f\\}(s):=\\int_{0}^{\\infty}e^{-st}f(t)dt.\n$$\n\n\\textbf{Lemma 2.3.}\\label{lemma} If $\\alpha\\leq{0}$ and $\\beta\\in\\mathbb{C}$, then the following recurrence formula holds:\n$$\nE_{\\alpha, \\beta}(z)=\\frac{1}{\\Gamma(\\beta)}+zE_{\\alpha, \\alpha+\\beta}(z).\n$$\nThis lemma was proved by R. K. Saxena in 2002 \\cite{L29}.\n\nLater, we use the properties of a Wright-type function studied by A. Pskhu \\cite{L30}, defined as\n$$\ne^{\\mu, \\delta}_{\\alpha, \\beta}(z)=\\sum_{n=0}^{\\infty}\\frac{z^n}{\\Gamma(\\alpha{n}+\\mu)\\Gamma(\\delta-\\beta{n})}, \\, \\, \\, \\alpha>0, \\, \\, \\alpha>\\beta.\n$$\n\nM-L function can be determined by Wright-type function as a special case $E_{\\alpha, \\beta}(z)=e^{\\beta, 1}_{\\alpha, 0}(z)$. So, we can record some properties of M-L function which can be reduced from the Wright-type function's properties.\n\n\\textbf{Lemma 2.4.} (\\cite{L30}) If $\\pi\\geq|argz|>\\frac{\\pi\\alpha}{2}+\\varepsilon, \\, \\, \\, \\varepsilon>0$, then the following relations are valid for $z \\to \\infty$:\n$$\n\\mathop {\\lim }\\limits_{\\mid{z}\\mid \\to \\infty} E_{\\alpha, \\beta}(z)=0,\n$$\n$$\n\\mathop {\\lim }\\limits_{\\mid{z}\\mid \\to \\infty} zE_{\\alpha, \\beta}(z)=-\\frac{1}{\\Gamma(\\beta-\\alpha)}.\n$$\n\n\n\n\\subsection{Regularized Caputo-like counterpart of the hyper-Bessel fractional differential operator}\n\n\\textbf{Definition 2.1.} (\\cite{L8}) The Riemann-Liouville fractional integral $I^{\\alpha}_{a+}f(t)$ and derivative $D^{\\alpha}_{a+}f(t)$ of order $\\alpha$ are defined by\n$$\nI^{\\alpha}_{a+}f(t)=\\frac{1}{\\Gamma(\\alpha)}\\int_{a}^{t}\\frac{f(\\tau)d\\tau}{(t-\\tau)^{1-\\alpha}},\n$$\n\n$$\nD^{\\alpha}_{a+}f(t)=\\left(\\frac{d}{dt}\\right)^nI^{n-\\alpha}_{a+}f(t), \\, \\, \\, \\, n-1\\leq\\alpha0, \\gamma\\in{\\mathbb{R}}$ and $\\beta>0$ is defined as (\\cite{L8})\n$$\nI^{\\gamma, \\delta}_{\\beta; a+}f(t)=\\frac{t^{-\\beta(\\gamma+\\delta)}}{\\Gamma(\\delta)}\\int_{a}^{t}(t^\\beta-\\tau^\\beta)^{\\delta-1}\\tau^{\\beta\\gamma}f(\\tau)d(\\tau^\\beta),\n$$\nwhich can be reduced up to a weight to $I^{q}_{a+}f(t)$ (Riemann-Liouville fractional integral) at $\\gamma=0$ and $\\beta=1$, and Erdelyi-Kober fractional derivative of $f(t)\\in{C}_{\\mu}^{(n)}$ for $n-1<\\delta\\leq{n}, n\\in\\mathbb{N}$, is defined by\n$$\nD_{\\beta, a+}^{\\gamma, \\delta}f(t)=\\prod_{j=1}^{n}\\left(\\gamma+j+\\frac{t}{\\beta}\\frac{d}{dt}\\right)\\big(I_{\\beta, a+}^{\\gamma+\\delta, n-\\delta} f(t)\\big),\n$$\nwhere $C_{\\mu}^{(n)}$ is the weighted space of continuous functions defined as\n$$\nC_{\\mu}^{(n)}=\\left\\{f(t)=t^p \\tilde{f(t)}; \\, \\, \\tilde{f} \\in{C}^{(n)}[0, \\infty)\\right\\}, \\, \\, \\, C_{\\mu}=C_{\\mu}^{(0)} \\, \\, \\textrm{with} \\, \\, \\mu\\in\\mathbb{R}.\n$$\n\n\\textbf{Definition 2.3.} Regularized Caputo-like counterpart of the hyper-Bessel fractional differential operator for $\\theta<1$, $0<\\alpha\\leq1$ and $t>a\\geq{0}$ is defined in terms of the E-K fractional order operator\n\n\\begin{equation}\\label{e2}\n^{C}\\Big(t^{\\theta}\\frac{d}{dt}\\Big)^\\alpha{f}(t)=(1-\\theta)^\\alpha{t}^{-\\alpha(1-\\theta)}D_{1-\\theta, a+}^{-\\alpha, \\alpha}\\left(f(t)-f(a)\\right)\n\\end{equation}\nor in terms of the hyper-Bessel differential (R-L type) operator\n\n\\begin{equation}\\label{eq3}\n^{C}\\left(t^{\\theta}\\frac{d}{dt}\\right)^\\alpha{f}(t)=\\Big(t^{\\theta}\\frac{d}{dt}\\Big)^\\alpha{f}(t)-\\frac{f(a)\\Big(t^{(1-\\theta)}-a^{(1-\\theta)}\\Big)^{-\\alpha}}{(1-\\theta)^{-\\alpha}\\Gamma(1-\\alpha)},\n\\end{equation}\nwhere\n$$\n\\left(t^\\theta\\frac{d}{dt}\\right)^\\alpha{f}(t)=\\left\\{ \\begin{gathered}\n(1-\\theta)^\\alpha t^{-(1-\\theta)\\alpha}I_{1-\\theta, a+}^{0, -\\alpha} f(t) \\, \\, \\, \\textrm{if} \\, \\, \\, \\theta<1, \\hfill \\cr\n(\\theta-1)^\\alpha t^{-(1-\\theta)\\alpha}I_{1-\\theta, a+}^{-1, -\\alpha} f(t) \\, \\, \\, \\textrm{if} \\, \\, \\, \\theta>1, \\hfill \\cr\n\\end{gathered} \\right.\n$$\nis a hyper-Bessel fractional differential operator (\\cite{L25}).\n\nFrom (\\ref{eq3}) for $a=0$ we obtain the definition presented in (\\cite{L25}) and also Caputo FDO is the particular case of Caputo-like counterpart hyper-Bessel operator at $\\theta=0$.\n\n\n\\textbf{Theorem 2.1.} Assume that the following conditions hold:\n\n$\\bullet$ $\\tau\\in{C}[0,1]$ such that $\\tau(0)=\\tau(1)=0$ and $\\tau'\\in{L}^2(0,1)$,\n\n$\\bullet$ $f(\\cdot, t)\\in{C}^3[0,1]$ and $f(x, \\cdot)\\in{C}_\\mu[a, T]$ such that\n\n $f(0,t)=f(\\pi,t)=f_{xx}(0,t)=f_{xx}(1,t)=0$, and $\\dfrac{\\partial^4}{\\partial{x}^4}f(\\cdot,t)\\in{L}^1(0,1).$\n\nThen, in $\\Omega=\\{0{0}.\n$$\n\nThis proves the uniqueness of the solution of the considered problem.\n\nMoreover, one can write an upper bound of $\\frac{1}{\\Delta}$ by using the last evaluation:\n\n\\begin{equation}\\label{e33}\n\t\\frac{1}{|\\Delta|}\\leq\\frac{M_1}{(k\\pi)^2}+M_2, \\, \\, \\, \\, (M_1, M_2=const).\n\\end{equation}\nNow we find an estimate for $B$ by using Lemma 2.1:\n\n$$\n|B|\\leq\\frac{\\lambda_k{p}^{1-\\alpha}}{\\Gamma(\\alpha)}|E_{\\delta, 1}(-\\lambda_k{a}^\\delta)|+\\lambda_k{a}^{\\delta-1}|E_{\\delta, \\delta}(-\\lambda_k{a}^\\delta)|\\leq\n$$\n$$\n\\leq\\frac{\\lambda_k{p}^{1-\\alpha}}{\\Gamma(\\alpha)}\\frac{M}{1+\\lambda_k{a}^\\delta}+\\lambda_k{a}^{\\delta-1}\\frac{M}{1+\\lambda_k{a}^\\delta}\\leq\n$$\n$$\n\\leq\\frac{\\lambda_k{p}^{1-\\alpha}}{\\Gamma(\\alpha)}\\frac{M}{\\lambda_k{a}^\\delta}+\\lambda_k{a}^{\\delta-1}\\frac{M}{\\lambda_k{a}^\\delta}=\\frac{Mp^{1-\\alpha}}{a^\\delta\\Gamma(\\alpha)}+\\frac{M}{a}=\\frac{M}{a}\\left(1+\\frac{p^{1-\\alpha}}{a^{\\delta-1}}\\right)=M_3, \\, \\, (M_3=const).\n$$\n\nUsing the last result and (\\ref{e33}) we estimate $\\left|\\frac{B}{\\Delta}\\varphi_k\\right|$:\n$$\n\\left|\\frac{B}{\\Delta}\\varphi_k\\right|\\leq\\frac{M_1M_3}{(k\\pi)^2}|\\varphi_k|+\\frac{M_2M_3}{k\\pi}|\\varphi_{1k}|\\leq\\frac{M_1M_3}{(k\\pi)^2}|\\varphi_k|+\\left(\\frac{M_2M_3}{k\\pi}\\right)^2+|\\varphi_{1k}|^2,\n$$\nhere $\\varphi_{1k}=2\\int_{0}^{1}\\varphi'(x)\\sin(k\\pi{x})dx$. Now let us find the upper bound of $C$:\n$$\n|C|\\leq\\int_{0}^{a}|a-s|^{\\delta+q-1}|E_{\\delta, \\delta+q}(-\\lambda_k(a-s)^\\delta)||f_k(s)|ds+\n$$\n$$\n+\\int_{0}^{a}|a-s|^{\\delta+q-2}|E_{\\delta, \\delta+q-1}(-\\lambda_k(a-s)^\\delta)||f_k(s)|ds\\leq\n$$\n$$\n\\leq\\int_{0}^{a}|a-s|^{\\delta+q-1}\\left|\\frac{M}{1+\\lambda_k|a-s|^\\delta}\\right||f_k(s)|ds+\n$$\n$$\n+\\int_{0}^{a}|a-s|^{\\delta+q-2}\\left|\\frac{M}{1+\\lambda_k|a-s|^\\delta}\\right||f_k(s)|ds\\leq\\frac{M_4}{(k\\pi)^2}, \\, \\, \\, (M_4=const.)\n$$\nHere we imply that $f_{x}(x,t)\\in{L}^1(0, a)$ for convergence of the last integral.\n\nThen, the estimate for $\\frac{C}{\\Delta}$ is\n\n$$\n\\left|\\frac{C}{\\Delta}\\right|\\leq\\frac{M_1M_4}{(k\\pi)^4}+\\frac{M_3M_4}{(k\\pi)^2}.\n$$\n\nFinally, we find the estimate for $|\\psi_k|$:\n\\begin{equation}\\label{e34}\n|\\psi_k|\\leq\\frac{M_1M_3}{(k\\pi)^2}|\\varphi_k|+\\left(\\frac{M_2M_3}{k\\pi}\\right)^2+|\\varphi_{1k}|^2+\\frac{M_1M_4}{(k\\pi)^4}+\\frac{M_3M_4}{(k\\pi)^2}<\\infty,\n\\end{equation}\nwhere $\\varphi'\\in{L}^2(0,1)$.\n\nFrom (\\ref{e32}) and in the same way one can show that $|\\tau_k|\\leq\\frac{M_5}{(k\\pi)^2}+|\\varphi_{1k}|^2<\\infty$, \\, \\, $M_5=const.$\n\n\\smallskip\n\nFor proving the existence of the solution, we need to show uniform convergency of series representations of $u(x, t)$, $u_{xx}(x,t)$, $^C\\left(t^\\theta\\frac{\\partial}{\\partial{t}}\\right)^{\\alpha}u(x, t)$ and $D_{t}^{(\\gamma, \\beta)\\mu}u(x, t)$ by using the solution (\\ref{e4}) and (\\ref{e27}) in $\\Omega_1$ and $\\Omega_2$ respectively.\n\n\nIn \\cite{L22}, the uniform convergence of series of $u(x,t)$ and $u_{xx}(x,t)$ showed for $t>0$. Similarly, for $t>a$, we obtain the following estimate:\n$$\n|u(x,t)|\\leq{M}\\sum_{k=1}^{\\infty}\\Big(\\frac{|\\tau_k|}{p^\\alpha+(k\\pi)^2|t^p-a^p|^\\alpha}+\\frac{1}{(k\\pi)^2}\\int_{a}^{t}|t^p-\\tau^p|^{\\alpha-1}f_{2k}(\\tau)d(\\tau^p)+\\\\\n$$\n$$\n+\\int_{a}^{t}\\frac{|t^p-\\tau^p|^{2\\alpha-1}}{p^\\alpha+(k\\pi)^2|t^p-a^p|^\\alpha}f_{2k}(\\tau)d(\\tau^p)\\Big),\n$$\nwhere $f_{2k}(t)=2\\int_{0}^{1}f_{xx}(x,t)\\sin(k\\pi{x})dx$.\n\nSince $|\\tau_k|<\\infty$ and $f(\\cdot, t)\\in{C}^3[0,1]$, then the above series converges and hence, by the Weierstrass M-test the series of $u(x, t)$ is uniformly convergent in $\\Omega_1$.\n\nThe series of $u_{xx}(x,t)$ is written in the form below\n$$\nu_{xx}(x,t)=-\\sum_{k=1}^{\\infty}(k\\pi)^2\\left(\\tau_kE_{\\alpha, 1}\\left[\\frac{(k\\pi)^2}{p^\\alpha}(t-a)^{p\\alpha}\\right]+G_k(t)\\right)\\sin(k\\pi{x}).\n$$\nWe obtain the following estimate:\n$$\n|u_{xx}(x,t)|\\leq{M}\\sum_{k=1}^{\\infty}\\Big(\\frac{(k\\pi)^2|\\tau_k|}{p^\\alpha+(k\\pi)^2|t^p-y^p|^\\alpha}+\\frac{1}{(k\\pi)^2}\\int_{a}^{t}|t^p-\\tau^p|^{\\alpha-1}|f_{4k}(\\tau)|d(\\tau^p)\\\\\n$$\n$$\n+\\int_{a}^{t}\\frac{|t^p-\\tau^p|^{2\\alpha-1}}{p^\\alpha+(k\\pi)^2|t^p-\\tau^p|^\\alpha}|f_{4k}(\\tau)|d(\\tau^p)\\Big),\n$$\nwhere $f_{4k}(t)=2\\int_{0}^{1}\\frac{\\partial^4}{\\partial{x}^4}f(x,t)\\sin(k\\pi{x})dx$ and $f(0,t)=f(1,t)=f_{xx}(0,t)=f_{xx}(1,t)=0$.\n\nSince $\\tau(0)=\\tau(1)=0$ and $\\frac{\\partial^4{f}}{\\partial{x}^4}(\\cdot, t)\\in{L}^1(0,1)$, then using integration by parts, we arrive at the following estimate\n\n$$\n|u_{xx}(x,t)|\\leq{M}\\sum_{k=1}^{\\infty}\\left(\\frac{1}{k}|\\tau_{1k}|+\\frac{1}{k^2}\\right)\\leq\\frac{M}{2}\\Big(\\sum_{k=1}^{\\infty}\\frac{3}{k^2}+\\sum_{k=1}^{\\infty}|\\tau_{1k}|^2\\Big),\n$$\nwhere $\\tau_{1k}=2\\int_{0}^{1}\\tau'(x)\\sin(k\\pi{x})dx$\nThen, the Bessel inequality for trigonometric functions implies\n$$\n|u_{xx}(x,t)|\\leq\\frac{M}{2}\\left(\\sum_{k=1}^{\\infty}\\frac{3}{k^2}+||\\tau'(x)||^2_{L^2(0,1)}\\right).\n$$\nThus, the series in the expression of $u_{xx}(x,t)$ is bounded by a convergent series which is uniformly convergent according to the Weierstrass M-test. Then, the series of $^C(t^\\theta\\frac{\\partial}{\\partial{t}})^\\alpha{u}(x,t)$ which can be written by\n$$\n^C\\left(t^\\theta\\frac{\\partial}{\\partial{t}}\\right)^\\alpha{u}(x,t)=-\\sum_{k=1}^{\\infty}(k\\pi)^2\\left(\\tau_kE_{\\alpha, 1}\\Big[-\\frac{(k\\pi)^2}{p^\\alpha}(t-a)^{p\\alpha}\\Big]+G_k(t)\\right)\\sin(k\\pi{x})+f(x,t),\n$$\nhas uniform convergence which can be showed in the same way to the uniform convergence of the series of $u_{xx}(x,t)$ (see \\cite{L22}).\n\nNow we need to show that the series of $u(x,t)$ and its derivatives should converge uniformly in $\\Omega_2$ by using (\\ref{e27}). We estimate\n\n$$\n|u(x,t)|\\leq|\\varphi_k||t^{(\\beta-2)(1-\\mu)}E_{\\delta, \\delta+\\mu(2-\\gamma)-1}(-\\lambda_k{t}^\\delta)|+|\\psi_k||t^{\\mu+(\\beta-1)(1-\\mu)}E_{\\delta, \\delta+\\mu(2-\\gamma)}(-\\lambda_k{t}^\\delta)|+\n$$\n$$\n+\\int_{0}^{1}|t-\\tau|^{\\delta-1}|E_{\\delta, \\delta}\\Big(-\\lambda_k(t-\\tau)^\\delta\\Big)||f_k(\\tau)|d\\tau.\n$$\nConsider estimates of the Mittag-Leffler function (see Lemma 2.1)\n$$\n|u(x,t)|\\leq\\frac{|t^{(\\beta-2)(1-\\mu)}||\\varphi_k|M}{1+\\lambda_k|{t}^\\delta)|}+\\frac{|t^{(\\beta-1)(1-\\mu)}||\\psi_k|M}{1+\\lambda_k|{t}^\\delta)|}\n$$\n$$\n+\\int_{0}^{t}|t-\\tau|^{\\delta-1}\\frac{M}{1+\\lambda_k|{(t-\\tau)}^\\delta)|}|f_k(\\tau)|d\\tau,\n$$\nwhere $f_{1k}(t)=\\int_{0}^{1}f_{x}(x,t)\\sin(k\\pi{x})dx$, $f_{x}(\\cdot,t)\\in{L^1[0,1]}$ and $\\varphi'\\in{L}^2[0,1]$. Then we obtain the estimate\n$$\n|u(x,t)|\\leq\\sum_{k=1}^{\\infty}\\frac{N_1}{(k\\pi)^2}, \\, \\, \\, \\, (N_1=\\textrm{const}),\n$$\nfor all $t>\\bar{t}>0, \\, \\, \\, 0\\leq{x}\\leq{1}$.\n\nIn the similar way one can show that\n$$\n|u_{xx}(x,t)|\\leq\\sum_{k=1}^{\\infty}(k\\pi)^2\\Big[|\\varphi_k||\\frac{M}{1+\\lambda_k{t}^\\delta}|+(\\frac{K_1}{(k\\pi)^2}+|\\varphi_{1k}|^2)|\\frac{M}{1+\\lambda_k{t}^\\delta}|+\n$$\n$$\n+\\int_{0}^{t}|t-\\tau|^{\\delta-1}\\frac{M}{1+\\lambda_k{t}^\\delta}|k_{2k}(\\tau)|d\\tau\\Big], \\, \\, \\, \\, (M=\\textrm{const}).\n$$\nThen, using Bessel's inequality and $\\varphi'\\in{L}^2[0,1]$, $f_{xxx}(\\cdot,t)\\in{L}^1(0,1)$, we get\n$$\n|u_{xx}(x,t)|\\leq\\frac{1}{2}\\Big(\\sum_{k=1}^{\\infty}\\frac{N_2}{(k\\pi)^2}+||\\varphi'(x)||^2_{L^2[0,1]}\\Big).\n$$\nWe have also used $2ab\\leq{a^2+b^2}.$\n\nUsing the equation in $\\Omega_2$, we write $D_{t}^{(\\gamma, \\beta)\\mu}u(x,t)$ in the form\n$$\nD_{t}^{(\\gamma, \\beta)\\mu}u(x,t)=u_{xx}(x,t)+f(x,t)$$\nand its uniform convergence can be done in a similar way to the uniform convergence of $u_{xx}(x,t)$ as\n$$\n|D^{(\\gamma, \\beta)\\mu}_tu(x,t)|\\leq\\sum_{k=1}^{\\infty}\\frac{N_3}{(k\\pi)^2}, \\, \\, \\, \\, (N_3=\\textrm{const}).\n$$\n\n\nFinally, considering the Weierstrass M-test, the above arguments prove that Fourier series in (\\ref{e4}) and (\\ref{e27}) converge uniformly in the domains $\\Omega_1$ and $\\Omega_2$. This is the proof that the considered problem's solution exists in $\\Omega$. Theorem 3.1 is proved.\n\n\n\n\\section*{Appendix}\n\nHere we write derivation of the series $^{C}\\left(t^{\\theta}\\frac{\\partial}{\\partial{t}}\\right)^\\alpha{u}(x,t)$ in (\\ref{e4}). Using relation (\\ref{eq3}) we get:\n$$\n^{C}\\left(t^{\\theta}\\frac{\\partial}{\\partial{t}}\\right)^\\alpha{u}(x,t)=\\sum_{k=0}^{\\infty}\\Big[(t^\\theta\\frac{\\partial}{\\partial{t}})^\\alpha\\left(\\tau_kE_{\\alpha,1}\\Big[-\\frac{(k\\pi)^2}{p^\\alpha}(t^p-a^p)^\\alpha\\Big]+G_k(t)\\right)\n-\\frac{\\tau_k(t^p-a^p)^\\alpha}{p^{-\\alpha}\\Gamma(1-\\alpha)}\\Big]\\sin(k\\pi{x}).\n$$\nThe hyper-Bessel derivative of the Mittag-Leffler function is\n$$\n\\Big(t^\\theta\\frac{\\partial}{\\partial{t}}\\Big)^\\alpha\\tau_kE_{\\alpha,1}\\left(-\\frac{(k\\pi)^2}{p^\\alpha}(t^p-a^p)^\\alpha\\right)=\\tau_kp^\\alpha(t^p-a^p)^{-\\alpha}E_{\\alpha, 1-\\alpha}\\big[-\\lambda(t^p-a^p)^\\alpha\\big].\n$$\nUsing the Lemma 2.3, we can write the last expression as follows\n$$\n\\Big(t^\\theta\\frac{\\partial}{\\partial{t}}\\Big)^\\alpha\\tau_kE_{\\alpha,1}\\Big[-\\frac{(k\\pi)^2}{p^\\alpha}(t^p-a^p)^\\alpha\\Big]=\\frac{\\tau_kp^\\alpha(t^p-a^p)^{-\\alpha}}{\\Gamma(1-\\alpha)}+\\tau_k(k\\pi)^2E_{\\alpha, 1}\\big[-\\frac{(k\\pi)^2}{p^\\alpha}(t^p-a^p)^\\alpha\\big].\n$$\n\nThen evaluating $\\Big(t^\\theta\\frac{\\partial}{\\partial{t}}\\Big)^\\alpha{G}_k(t)$ gives that\n$$\n\\Big(t^\\theta\\frac{\\partial}{\\partial{t}}\\Big)^\\alpha{G}_k(t)=\\Big(t^\\theta\\frac{\\partial}{\\partial{t}}\\Big)^\\alpha\\Big(f^*_k(t)+\\lambda^*\\int_{a}^{t}(t^p-a^p)^{\\alpha-1}E_{\\alpha, \\alpha}\\big[\\lambda^*(t^p-a^p)\\big]f^*_k(\\tau)d(\\tau^p)\\Big)=\n$$\n$$\n=p^\\alpha{t}^{-p\\alpha}D_{p, a+}^{-\\alpha, \\alpha}\\Big(\\frac{1}{p^\\alpha}I_{p, a+}^{-\\alpha, \\alpha}t^p\\alpha{f}_k(t)+\\lambda^*\\int_{a}^{t}(t^p-a^p)^{\\alpha-1}E_{\\alpha, \\alpha}\\big[\\lambda^*(t^p-a^p)\\big]f^*_k(\\tau)d(\\tau^p)\\Big)=\n$$\n$$\n=f_{k}(t)+p^{\\alpha}t^{-p\\alpha}D_{p, a+}^{-\\alpha, \\alpha}\\Big(\\lambda^*\\int_{a}^{t}(t^p-a^p)^{\\alpha-1}E_{\\alpha, \\alpha}\\big[\\lambda^*(t^p-a^p)\\big]f^*_k(\\tau)d(\\tau^p)\\Big),\n$$\nwhere $\\lambda^*=-\\frac{\\lambda_k}{p^\\alpha}$ and $f^*_k(t)=\\frac{1}{p^\\alpha\\Gamma(\\alpha)}\\int_{a}^{t}(t^p-\\tau^p)^{\\alpha-1}f_k(\\tau)d(\\tau^p)$.\n\nThe second term in the last expression can be simplified using the Erd'elyi-Kober fractional derivative for $n=1$,\n$$\n-\\lambda_kt^{-p\\alpha}\\left(1-\\alpha+\\frac{t}{p}\\frac{d}{dt}\\right)\\frac{t^{-p(1-\\alpha)}}{\\Gamma(1-\\alpha)}\\int_{a}^{t}(t^p-\\tau^p)^{-\\alpha}d(\\tau^p)\\int_{a}^{\\tau}(\\tau^p-s^p)^{\\alpha-1}E_{\\alpha, \\alpha}\\big[\\lambda^*(\\tau^p-s^p)^\\alpha\\big]f^*_k(s)d(s^p)=\n$$\n$$\n-\\lambda_kt^{-p\\alpha}\\left(1-\\alpha+\\frac{t}{p}\\frac{d}{dt}\\right)\\frac{t^{-p(1-\\alpha)}}{\\Gamma(1-\\alpha)}\\int_{a}^{t}f^*_k(s)d(s^p)\\int_{s}^{t}(t^p-\\tau^p)^{-\\alpha}(\\tau^p-s^p)^{\\alpha-1}E_{\\alpha, \\alpha}\\big[\\lambda^*(\\tau^p-s^p)^\\alpha\\big]d(\\tau^p)=\n$$\n$$\n-\\lambda_kt^{-p\\alpha}\\left(1-\\alpha+\\frac{t}{p}\\frac{d}{dt}\\right)t^{-p(1-\\alpha)}\\int_{a}^{t}E_{\\alpha, 1}\\left[\\lambda^*(t^p-s^p)^{\\alpha}\\right]f^*_k(s)d(s^p)=\n$$\n$$\n-\\lambda_k(1-\\alpha)t^{-p}\\int_{a}^{t}E_{\\alpha, 1}\\left[\\lambda^*(t^p-s^p)^{\\alpha}\\right]f^*_k(s)d(s^p)-\n$$\n$$\n-\\frac{\\lambda_kt^{-p\\alpha+1}}{p}\\frac{d}{dt}\\left(t^{-p(1-\\alpha)}\\int_{a}^{t}E_{\\alpha, 1}\\left[\\lambda^*(t^p-s^p)^{\\alpha}\\right]f^*_k(s)d(s^p)\\right)=\n$$\n$$\n=-\\lambda_k(1-\\alpha)t^{-p}\\int_{a}^{t}E_{\\alpha, 1}\\left[\\lambda^*(t^p-s^p)^{\\alpha}\\right]f^*_k(s)d(s^p)+\n$$\n$$\n+\\lambda_k(1-\\alpha)t^{-p}\\int_{a}^{t}E_{\\alpha, 1}\\left[\\lambda^*(t^p-\\tau^p)^{\\alpha}\\right]f^*_k(\\tau)d(\\tau^p)-\\lambda_kf^*_k(t)-\n$$\n$$\n-\\lambda_kt^{1-p}\\int_{a}^{t}\\lambda^*(t^p-\\tau^p)^{\\alpha-1}E_{\\alpha, \\alpha}\\left[\\lambda^*(t^p-\\tau^p)^{\\alpha}\\right]f^*_k(\\tau)d(\\tau^p)=\n$$\n$$\n=-\\lambda_k\\left(f^*_k(t)+\\lambda^*\\int_{a}^{t}(t^p-a^p)^{\\alpha-1}E_{\\alpha, \\alpha}\\big[\\lambda^*(t^p-a^p)\\big]f^*_k(\\tau)d(\\tau^p)\\right)=-\\lambda_kG_k(t).\n$$\nHence, we get\n$$\n^{C}\\left(t^{\\theta}\\frac{\\partial}{\\partial{t}}\\right)^\\alpha{u}(x,t)=-\\sum_{k=0}^{\\infty}(k\\pi)^2\\left[\\tau_kE_{\\alpha, 1}\\left(\\frac{(k\\pi)^2}{p^\\alpha}(t^p-a^p)\\right)+G_k(t)\\right]\\sin(k\\pi{x})+f(x,t).\n$$\nThis proves that solution (\\ref{e4}) satisfies the equation\n$$\n^C\\Big(t^\\theta\\frac{\\partial}{\\partial{t}}\\Big)^{\\alpha}u(x,t)-u_{xx}(x,t)=f(x,t).\n$$\n\nWe would like to note that using the result of this work, one can consider FPDE with the Bessel operator considering local \\cite{krma} and non-local boundary value problems \\cite{karruzh}. In that case the Fourier-Bessel series will play an important role. The other possible applications are related with the consideration of more general operators in space variables. For instance, in \\cite{rtb1}, very general positive operators have been considered, and the results of this paper can be extended to that setting as well.\n\n\\section{Acknowledgement}\n\n\nThe second author was partially supported by the FWO Odysseus 1 grant G.0H94.18N: Analysis and Partial Differential Equations and by the Methusalem programme of the Ghent University Special Research Fund (BOF) (Grant number 01M01021). \n\n\\bigskip\n\n\n\\textbf{\\Large References}\n\n\\begin{enumerate}\n\\bibitem{L1} {\\it C.~Friedrich.} Relaxation and retardation functions of the Maxwell model with fractional derivatives. Rheologica Acta. 30(2), 1991, pp. 151-158.\n\n\\bibitem{L2} {\\it D.~Kumar, J.~Singh , M.~Al Qurashi.} A new fractional SIRS-SI malaria disease model with application of vaccines, antimalarial drugs,\nand spraying. Adv Differ Equa. (2019), 2019: 278.\n\n\\bibitem{L3} {\\it I.~Podlubny.} Fractional Differential Equations. United States, Academic Press.\n1999. 340 p.\n\n\\bibitem{L4} {\\it D.~Baleanu, Z.~B.~Guvenc and J.~A.~T.~Machado.} New trends in nanotechnology and fractional calculus applications, Computers and Mathematics with Applications. 59, 2010, pp. 1835-1841.\n\n\\bibitem{L5} {\\it C.~G.~Koh and J.~M.~Kelly.} Application for a fractional derivative to seismic analysis of base-isolated models, Earthquake Engineering and Structural Dynamics. 19, 1990, pp. 229-241.\n\n\\bibitem{L6} {\\it E.~T.~Karimov, A.~S.~Berdyshev, N.~A.~Rakhmatullaeva.} Unique solvability of a\nnon-local problem for mixed-type equation with fractional derivative. {Mathematical\n\tMethods in the Applied Sciences.} 40(8), 2017, pp. 2994-2999.\n\n\\bibitem{L7} {\\it Z.~A.~Nakhusheva.} Non-local boundary value problems for main and mixed type differential equations. Nalchik, 2011 [in Russian].\n\n\\bibitem{L8} {\\it A.~A.~Kilbas, H.~M.~Srivastava, J.~J.~Trujillo.} Theory and\nApplications of Fractional Differential Equation. Elsevier, Amsterdam. 2006.\n\n\n\\bibitem{sad} {\\it M.~Kirane, M.~A.~Sadybekov, A.~A.~Sarsenbi.} On an inverse problem of reconstructing a subdiffusion process from nonlocal data. {Mathematical\n\tMethods in the Applied Sciences.} 42(6), 2019, pp.2043-2052.\n\n\\bibitem{ruzh} {\\it M.~Ruzhansky, N.~Tokmagambetov and B.~Torebek.} On a non-local problem for a multi-term fractional diffusion-wave equation. Fractional Calculus and Applied Analysis, 23(2), 2020, pp. 324-355.\n\n\\bibitem{karruzh} {\\it E.~Karimov, M.~Mamchuev and M.~Ruzhansky.} Non-local initial problem for second order time-fractional and space-singular equation. Hokkaido Math. J., 49, 2020, pp.349-361.\n\n\\bibitem{L9} {\\it E.~T.~Karimov.} Boundary value problems for parabolic-hyperbolic type equations with spectral para\\-meters. PhD Thesis, Tashkent, 2006.\n\n\\bibitem{L10} {\\it S.~Kh.~Gekkieva.} A boundary value problem for the generalized transfer equation with a fractional derivative in a semi-infinite domain. Izv. Kabardino-Balkarsk. Nauchnogo Tsentra RAN, 8(1), 2002, pp. 3268-3273. [in Russian].\n\n\\bibitem{L11} {\\it E.~T.~Karimov, J.~S.~Akhatov.} A boundary problem with integral gluing\ncondition for a parabolic-hyperbolic equation in volving the Caputo fractional derivative.\n{Electronic Journal of Differential Equations.} 14, 2014, pp. 1-6.\n\\bibitem{L12} {\\it P. ~Agarwal, A.~Berdyshev and E.~Karimov.} Solvability of a non-local\nproblem with integral form transmitting condition for mixed type equation with Caputo\nfractional derivative. {Results in Mathematics.} 71(3), 2017, pp.1235-1257.\n\\bibitem{L13} {\\it B.~J.~Kadirkulov.} Boundary problems for mixed parabolic-hyperbolic\nequations with two lines of changing type and fractional derivative. {Electronic Journal\n\tof Differential Equations}, 2014(57), pp. 1-7.\n\\bibitem{L14} {\\it R.~Hilfer.} Fractional time evolution, in: R. Hilfer (ed.), Applications of Fractional Calculus in Physics, World SCi., Singapore, 2000, pp. 87-130.\n\\bibitem{L15} {\\it R.~Hilfer.} Experimental evidence for fractional time evolution in glass forming materials. J. Chem. Phys. 284, 2002, pp. 399-408.\n\\bibitem{L16} {\\it R.~Hilfer, Y.~Luchko and Z.~Tomovski.} Operational method for solution of the fractional differential equations with the generalized Riemann-Liouville fractional derivatives. Fractional Calculus and Applied Analysis. 12, 2009, pp. 299-318.\n\n\\bibitem{L17} {\\it O.~Kh.~Abdullaev, K.~Sadarangani.} Non-local problems with integral gluing\ncondition for loaded mixed type equations involving the Caputo fractional derivative.\n\\textit{Electronic Journal of Differential Equations}, 164, 2016, pp. 1-10.\n\n\\bibitem{L18} {\\it A.~S.~Berdyshev, B.~E.~Eshmatov, B.~J.~Kadirkulov.} Boundary value problems\nfor fourth-order mixed type equation with fractional derivative.\n{Electronic Journal of Differential Equations}, 36, 2016, pp. 1-11.\n\\bibitem{L19} {\\it E.~T.~Karimov.} Tricomi type boundary value problem with integral conjugation\ncondition for a mixed type equation with Hilfer fractional operator. {Bulletin of the\n\tInstitute of Mathematics}. 1, 2019, pp. 19-26.\n\\bibitem{L20} {\\it V.~M.~Bulavatsky.} Closed form of the solutions of some boundary-value problems for anomalous diffusion equation with Hilfer's generalized derivative. Cybernetics and Systems Analysis. 30(4), 2014, pp. 570-577.\n\n\\bibitem{Ldn} {\\it M.~M.~Dzhrbashyan, A.~B.~Nersesyan.} Fractional Derivatives and the\nCauchy Problem for Fractional Differential Equations, Izv. Akad. Nauk\nArmyan. SSR. 3, No 1 (1968), 3\u201329.\n\\bibitem{bag} {\\it F.~T.~Bogatyreva.} Initial value problem for fractional order equation with \tconstant coefficients, Vestnik KRAUNC. Fiz.-mat. nauki. 2016, 16: 4-1, 21-26. DOI: 10.18454\/2079- 6641-2016-16-4-1-21-26\n\\bibitem{bag2} {\\it F.~T.~Bogatyreva.} Representation of solution for first-order partial differential equation\nwith Dzhrbashyan \u2013 Nersesyan operator of fractional differentiation. Reports ADYGE (Circassian) International Academy of Sciences. Volume 20. \u2116 2. pp.6-11. \n\\bibitem{Lfca} {\\it M.~M.~Dzherbashian, A.~B.~Nersesian.} Fractional derivatives\n\tand Cauchy problem for differential equations of fractional order. Fract. Calc. Appl. Anal. 23, No 6 (2020), 1810-1836.\tDOI: 10.1515\/fca-2020-0090\n\n \\bibitem{L21} {\\it I.~Dimovski.} Operational calculus for a class of differential operators, C.R. Acad. Bulg. Sci. 19(12), 1996, pp.1111-1114.\n\n\\bibitem{L22} {\\it F.~Al-Musalhi, N.~Al-Salti and E.~Karimov.} Initial boundary value\nproblems for fractional differential equation with hyper-Bessel operator.\nFractional Calculus and Applied Analysis. 21(1), 2018, pp. 200-219.\n\n\\bibitem{L23} {\\it N.~H.~Tuan, L.~N.~Huynh, D.~Baleanu, N.~H.~Can.} On a terminal value problem for a\n\tgeneralization of the fractional diffusion equation with hyper-Bessel operator. {Mathematical\n\tMethods in the Applied Sciences.} 43(6), 2019, pp.2858-2882.\n\n\\bibitem{L24} {\\it K.~Zhang.} Positive solution of nonlinear fractional differential equations with\n\tCaputo-like counterpart hyper-Bessel operators. {Mathematical\n\tMethods in the Applied Sciences.} 43(6), 2019, pp. 2845-2857.\n\n\\bibitem{L25} {\\it R.~Garra, A.~Giusti, F.~Mainardi, G.~Pagnini.} Fractional relaxation\nwith time-varying coefficient. Fractional Calculus and Applied Analysis. 17(2), 2014, pp. 424-439.\n\n\\bibitem{L26} {\\it E.~T.~Karimov, B.~H.~Toshtemirov.} Tricomi type problem with integral conjugation condition for a mixed type equation with the hyper-Bessel fractional differential operator. Bulletin of the Institute of Mathematics. 4, 2019, pp. 9-14.\n\n\\bibitem{L27} {\\it F.~Mainardi.} On some properties of the Mittag-Leffler function $E_{\\alpha}(-t^{\\alpha})$, completely monotone for $t>0$ with $0<\\alpha<1$. Discrete and Continuous Dynamical Systems - B. 19 (7), 2014, pp. 2267-2278.\n\n\\bibitem{L28} {\\it L.~Boudabsa, T.~Simon.} Some Properties of the Kilbas-Saigo Function. Mathematics. 9(3), 2021: 217.\n\n\\bibitem{L29} {\\it R.~K.~Saxena.} Certain properties of generalized Mittag-Leffler function, in Proceedings of the 3rd\nAnnual Conference of the Society for Special Functions and Their Applications, pp. 77-81, Chennai, India, 2002.\n\\bibitem{L30} {\\it A.~V.~Pskhu.} Partial Differential Equations of Fractional Order. Moscow, Nauka. 2005. [in Russian].\n\n\\bibitem{KDSc} {\\it E.~T.~Karimov.} Boundary value problems with integral transmitting conditions and inverse problems for integer and fractional order differential equations. DSc Thesis, Tashkent, 2020.\n\n\\bibitem{L31} {\\it E.~I.~Moiseev.} On the basis property of systems of sines and cosines.\nDoklady AN SSSR. 275(4), 1984, pp.794-798.\n\n\\bibitem{krma} {\\it P.~Agarwal, E.~Karimov, M.~Mamchuev, M.~ Ruzhansky.} On boundary-value problems for a partial differential equation with Caputo and Bessel operators, in Recent Applications of Harmonic Analysis to Function Spaces, Differential Equations, and Data Science: Novel Methods in Harmonic Analysis, Vol 2, Appl. Numer. Harmon. Anal., 707-718, Birkhauser\/Springer, 2017.\n\n\\bibitem{rtb1} {\\it M.~ Ruzhansky, N.~Tokmagambetov, B.~Torebek.} Inverse source problems for positive operators. I: Hypoelliptic diffusion and subdiffusion equations. Journal of Inverse and Ill-Posed Problems, 27 (2019), 891-911.\n\\end{enumerate}\n\n\n\n\\end{document}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzbxze b/data_all_eng_slimpj/shuffled/split2/finalzzbxze new file mode 100644 index 0000000000000000000000000000000000000000..d32f1217a33f81d93f501786a395092829592dc3 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzbxze @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\nOptimal investment with intermediate consumption and a stream of labor income (or liabilities) is one of the central problems in mathematical economics. \nIf borrowing against the future income is prohibited, the main technical difficulty lies in the fact that there are infinitely many constraints. Even in the deterministic case of no stocks and non-random income, a classical approach is based on the {\\it convexification of the constraints} that leads to a non-trivial dual problem formulated over decreasing nonnegative functions.\n\nBorrowing constraints imposed at all times not only affect the notion of admissibility, leading to more difficult mathematical analysis, but also\nchange the meaning to fundamental concepts of mathematical finance such as replicability and completeness. The latter is formulated via the attainability of every (bounded) contingent claim by a portfolio of traded assets. For a labor income\/liability streams that pays off dynamically, \nthere is no a priori guarantee that such a replicating portfolio (if it exists at all) is admissible, i.e., satisfies the constraints. Thus, in the terminology of \\cite{HePages}, even a complete market becomes dynamically incomplete under the borrowing constraints. \n The analysis of such a problem (in otherwise complete Brownian settings with a corresponding unique risk-neutral measure), is performed in \\cite{HePages} and later in \\cite{ElKarouiJeanblanc}. \nThe nonnegative decreasing {\\it processes} (that parametrize the dynamic incompleteness mentioned above) play an important role in the characterizations of optimal investment and consumption plans. The analysis in \\cite{HePages} and \\cite{ElKarouiJeanblanc} is connected with optimal stopping techniques from~\\cite{Kar89}. \n\nIncomplete markets with no-borrowing constraints have been analyzed only in specific Markovian models in \\cite{DuffieZarip1993} and \\cite{DuffieFlemingSonerZarip1997} based on partial differential equations techniques. \n The goal of the present paper is to study the problem of consumption and investment with no-borrowing constraints in general (so, non-Markovian) incomplete models. This leads to having, simultaneously, two layers of incompleteness. One comes from the many martingale measures, the other from a similar class of non-decreasing processes (as above) that describe the dynamic incompleteness. \n We refer to \\cite{MalTrub07} for the examples of market incompleteness in finance and macroeconomics. \n\nIn contrast to \\cite{HePages} and \\cite{ElKarouiJeanblanc}, our model not only allows for incompleteness, but also for jumps. Mathematically, this means we choose to work in a general semimartingale framework.\n\nAs in \\cite{HePages} and \\cite{ElKarouiJeanblanc}, \nour approach is based on duality. One of the principal difficulties is in the construction of the dual feasible set and the dual value function. It is well-known that the martingale measures drive the dual domain in many problems of mathematical finance. On the other hand, the convexification of constraints leads to the decreasing processes as the central dual object as well. We show that the dual elements in the incomplete case can be approximated by products of the densities of martingale measures and such nonnegative decreasing processes. This is one of the primary results of this work, see section \\ref{secStructureOfDualDomain}, that leads to the complementary slackness characterization of optimal wealth in section \\ref{secSlackness}, where it is shown that the approximating sequence for the dual minimizer leads to a nonincreasing process, which decreases at most when the constraints are attained. In turn, the dual minimizer can be written as a product of such a nondecreasing process and an optional strong supermartingale deflator. In the case of complete Brownian markets a similar result is proved in \\cite{HePages} and \\cite{ElKarouiJeanblanc}.\n\nIn order to implement the approach, we {\\it increase the dimensionality} of the problem and treat as arguments of the indirect utility not only the initial wealth, but also the function that specifies the number of units of labor income (or the stream of liabilities) at any later time. This parametrization has the spirit of \\cite{HK04}, however unlike \\cite{HK04}, we go into (infinite-dimensional) non-reflexive spaces, which gives both novelty and technical difficulties to our analysis.\n Also, our formulation permits to price by marginal rate of substitution the whole labor income process. This is done \n\nthrough the subdifferentiability results in section \\ref{secSubdifferentiability}. Note that the sub differential elements (prices) are time-dependent, so infinite-dimensional, unlike in \\cite{HK04}.\n\n\nAnother contribution of the paper lies in the {\\it unified framework of admissibility} outlined in section \\ref{secModel}. More precisely, we assume that no-borrowing constraints are imposed starting from some pre-specified stopping time and hold up to the terminal time horizon. \nThis framework allows us to treat in one formulation both the problem of no-borrowing constraints at all times (described above) and the one where borrowing against the future income is permitted with a constraint only at the end. The latter is well-studied in the literature, see \\cite{CvitanicSchachermayerWang}, \\cite{Karatzas-Zitkovic-2003}, \\cite{HK04}, \\cite{Zitkovic}, and \\cite{MostovyiRandEnd}.\n In such a formulation, the constraints reduce to a single inequality and the decreasing processes in the dual feasible set become constants. \n \nAmong the many possibilities of constraints, \\cite{Cuoco1997} considers the problem of investment and consumption with labor income and no-borrowing constraints in Brownian market, even allowing for incompleteness. The dual problem cannot be solved directly (in part because for these constraints the dual space considered is too small) but the primal can be solved with direct methods. An approximate dual sequence can then be recovered from the primal. We generalize \\cite{HePages} and \\cite{ElKarouiJeanblanc} (complete Brownian markets) and \\cite{Cuoco1997} (possibly incomplete Brownian markets) to the case of general semi-martingale incomplete markets. Our dual approach allows us, at the same time to obtain a dual characterization (complementary slackness) of the optimal consumption plan (not present in \\cite{Cuoco1997}), similar to the complete case in \\cite{HePages} and \\cite{ElKarouiJeanblanc} and the possibility to study the dependence on labor income streams, through the parametrization of such streams.\n \n\n\n\nEmbedding path dependent problems into the convex duality framework have been analyzed in \\cite{XiangYuHabit}, \\cite{HaoMat}, \\cite[Section 3.3]{Pham01}, whereas without duality but with random endowment it is considered in \\cite{Miklos2018}, in the abstract singular control setting the duality approach is investigated in \\cite{BankKaupila}.\nOur embedding does not require any condition on labor income {replicability}, which becomes highly technical in the presence of extra admissibility constraints. Even in the case where the only constraint is imposed at maturity, in this part our approach differs from the one in \\cite{HK04}, where non-replicability of the endowment (in the appropriate sense) is used in the proofs as it ensures that the effective domain of the dual problem has the same dimensionality as the primal domain. Note that even if \n the labor income is spanned by the same sources of randomness as the stocks, the idea of replicating the labor income and then reducing the problem to the one without it, does not necessarily work under the borrowing constraints, see the discussion in \\cite[pp. 671-673]{HePages}. \n\n\n\n\n\n\n\n Some of the more specific technical contributions of this paper can be summarized as follows:\n\n\\begin{itemize}\n\\item We analyze the {\\it boundary behavior} of the value functions. Note that the value functions are defined over infinite-dimensional spaces. \n\n\n\\item The finiteness of the indirect utilities {\\it without} labor income is imposed only, as a necessary and sufficient condition that allows for the standard conclusions of the utility maximization theory, see \\cite{MostovyiNec}.\n\n\n\\item We show existence-uniqueness results for the {\\it unbounded} labor income both from above and below.\n\n\n\n\\item We observe that the ``Snell envelope proposition'' \\cite[Proposition 4.3]{K96} can be extended to \n the envelope over all stopping times that exceed a given initial stopping time $\\theta _0$.\n\\item We represent the dual value function in terms of {\\it uniformly integrable} densities of martingale measures, i.e., the densities of martingale measures under which the maximal wealth process of a self-financing portfolio that superreplicates the labor income, is a uniformly integrable martingale, see Lemma \\ref{finitenessOverZ'} below. \n\n\\end{itemize}\n\n\n\n{\\bf Organization of the paper}. In Section \\ref{secModel}, we specify the model. We state and prove existence, uniqueness, semicontinuity and biconjugacy results in Section \\ref{secProofs}, subdifferentiability is proven in Section \\ref{secSubdifferentiability}. Structure of the dual domain is analyzed in Section \\ref{secStructureOfDualDomain} and complimentary slackness is established in Section \\ref{secSlackness}.\n\\section{Model}\\label{secModel}\n\nWe consider a {financial market model} with finite time horizon\n$[0,T]$ and a zero interest rate. The price process $S=(S^i)_{i=1}^d$\nof the stocks is assumed to be a\nsemimartingale on a complete stochastic basis\n$\\left( \\Omega, \\mathcal F, \\left(\\mathcal F_t \\right)_{t\\in[0,T]}, \\mathbb\nP\\right)$, where $\\mathcal F_0$ is trivial. \n\nLet $(e)_{t\\in[0,T]}$ be an optional process that specifies the labor income rate, which is assumed to follow a certain stochastic clock, that we specify below.\nBoth processes $S$ and $e$ are given exogenously. \n\nWe define a \\textit{stochastic clock} as a\nnondecreasing, c\\`adl\\`ag, adapted process such that\n\\begin{equation}\n \\tag{finClock}\n \\label{finClock}\n\\kappa_0 = 0, ~~ \\mathbb P\\left[\\kappa_T>0 \\right]>0,\\text{ and } \\kappa_T\\leq A\n\\end{equation}\nfor some finite constant $A$. We note that the stochastic clock allows to include multiple standard formulations of the utility maximization problem in one formulation, see e.g., \\cite[Example 2.5 - 2.9]{MostovyiNec}.\nLet us define \n\\begin{equation}\\label{defK}\nK_t {:=} \\mathbb E\\left[\\kappa_t \\right],\\quad t\\in[0,T].\n\\end{equation}\n\\begin{Remark}\nThe function $K$ defined in \\eqref{defK} is right-continuous with left limits and takes values in $[0,A]$. \n\\end{Remark}\nWe assume the income and consumption are given in terms of the clock $\\kappa$.\nDefine a {portfolio} $\\Pi$ as a quadruple $(x, q, H, c),$ where the\nconstant $x$ is the initial value of the portfolio, the function $q:[0,T]\\to\\mathbb R$ is a bounded and Borel measurable function,\n which \nspecifies the amount of labor income rate, $H =\n(H_i)_{i=1}^d$ is a predictable $S$-integrable process that\ncorresponds to the amount of each stock in the portfolio, and\n$c=(c_t)_{t\\in[0,T]}$ is the consumption rate, which we assume to be\n optional and nonnegative.\n\nThe \\textit{wealth process} $V=(V_t)_{t\\in[0,T]}$ {generated by the}\nportfolio is \n\\begin{equation}\\nonumber\nV_t = x + \\int_0^t H_sdS_s + \\int_0^t\\left(q_se_s-c_s\\right)d\\kappa_s,\\quad t\\in[0,T].\n\\end{equation}\nA portfolio $\\Pi$ with $c\\equiv 0$ and $q \\equiv 0$ is called\n\\textit{self-financing}. The collection of nonnegative wealth\nprocesses {generated by} self-financing portfolios with initial value $x\\geq 0$ is\ndenoted by $\\mathcal X(x)$, i.e.\n\\begin{equation}\\nonumbe\n\\mathcal X(x) {:=} \\left\\{X\\geq 0:~X_t=x + \\int_0^t\nH_sdS_s,~~t\\in[0,T] \\right\\},~~x\\geq 0.\n\\end{equation}\nA probability measure $\\mathbb Q$ is an \\textit{equivalent local\nmartingale measure} if $\\mathbb Q$ is equivalent to $\\mathbb P$ and\nevery $X\\in\\mathcal X(1)$ is a local martingale under $\\mathbb Q$. We\ndenote the family of equivalent local martingale measures by\n$\\mathcal M$ and assume that\n\\begin{equation}\\tag{noArb}\\label{NFLVR}\n\\mathcal M \\neq \\emptyset.\n\\end{equation}\nThis condition is equivalent to the absence of arbitrage\nopportunities \n{in} the marke\n, see\n \\cite{DS, DS1998} as well as \n\\cite{KarKar07} for the exact statements and further\nreferences.\n\nTo rule out doubling strategies in the presence of random endowment,\nwe need to impose additional restrictions. Following \\cite{DS97}, we say that a\nnonnegative process in $\\mathcal X(x)$ is \\textit{maximal} if its\nterminal value cannot be dominated by that of any other process in\n$\\mathcal X(x)$.\nAs in \\cite{DS97}, we define an\n\\textit{acceptable} process to be a process of the form $X = X' -\nX'',$ where $X'$ is a nonnegative wealth process {generated by} a self-financing portfolio and $X''$ is maximal.\n\n Our unified framework of admissibility is given by a fixed stopping time $\\theta_0$. The no-borrowing constraints will hold starting at this stopping time until the end. Let $\\Theta$ be the set of stopping times that are greater or equal than $\\theta_0$. \n\\begin{Lemma}\\label{3311}\nLet $q^1$ and $q^2$ be bounded, Borel measurable functions on $[0,T]$, such that $q^1= q^2$, $dK$-a.e. Then, the cumulative labor income processes $\\int_0^{\\cdot}q^i_se_sd\\kappa_s$, $i = 1,2$, are indistinguishable. \n\\end{Lemma} The proof of this lemma is given in section \\ref{subSecCharDualDomain}.\nFollowing\n \\cite{HK04}, we denote by\n$\\mathcal X(x,q)$ the set of acceptable processes with initial {values}\n$x$, that dominate the labor income on $\\Theta$:\n\\begin{equation}\\nonumbe\n\\begin{array}{rcl}\n\\mathcal X(x, q) &{:=}& \\left\\{ \n{\\rm acceptable}~X: X_0 = x~{\\rm and}\\right.\\\\\n&&\\left.~ X_{\\tau} + \\int_0^{\\tau}q_se_sd\\kappa_s \\geq 0, ~\\mathbb P-a.s.~{\\rm for~every~}\\tau\\in\\Theta\n\\right\\}.\\\\\n\\end{array}\n\\end{equation}\nLet us set\n\\begin{equation}\\nonumbe\n\\mathcal K {:=} \\left\\{(x,q): \\mathcal X(x,q) \\neq \\emptyset \\right\\},\n\\end{equation}\nLet $\\mathring{\\mathcal K}$ denote the interior or $\\mathcal K$ in the $\\mathbb R\\times \\mathbb L^\\infty(dK)$-norm topology.\nWe characterize $\\mathcal K$ in Lemma~\\ref{lemma1} below that in particular asserts that under Assumption \\ref{asEnd}, $\\mathring{\\mathcal K}\\neq \\emptyset$. \n \nThe set of admissible consumptions is defined as\n\\begin{equation}\\nonumbe\n\\begin{array}{rcl}\n\\mathcal A(x,q)&{:=} &\\left\\{ \noptional ~c\\geq 0:~~there~exists~X\\in\\mathcal X(x,q), ~such~that~ \\right.\\\\\n&& \n\\left.\\int_0^{\\tau}c_sd\\kappa_s\\leq X_{\\tau} + \\int_0^{\\tau}q_se_sd\\kappa_s,\n~~for~every~\\tau\\in\\Theta\\right\\},\\quad (x,q)\\in\\mathcal K.\\\\ \n\\end{array}\n\\end{equation}\nNote that $c\\equiv 0$ belongs to $\\mathcal A(x,q)$ for every $(x,q)\\in\\mathcal K$. \n\\begin{Remark}\\label{Remark-Mihai} The no-borrowing constraints can also be written as\n$$\n\\mathbb{P}\\left (\\int_0^{t}c_sd\\kappa_s\\leq X_{t} + \\int_0^{t}q_se_sd\\kappa_s,\n~~for~every~ \\theta _0\\leq t\\leq T \\right)=1.$$\nWe write the constraints in terms of stopping times $\\tau \\in \\Theta$ as \nwe use the stopping times $\\tau \\in \\Theta$ (and the corresponding decreasing processes that jump from one to zero at these times) as the building blocks of our analysis. \n\\end{Remark}\n\\begin{Remark}\\label{3312}\nIt follows from Lemma \\ref{3311}, for every $x\\in\\mathbb R$, we have\n$$\\mathcal X(x,q^1) = \\mathcal X(x,q^2)\\quad and \\quad \\mathcal A(x,q^1) = \\mathcal A(x, q^2),$$\nwhere some of these sets might be empty\n\\end{Remark} \n\nHereafter, we shall impose the following conditions on the endowment process.\n\\begin{Assumption}\\label{asEnd}\nThere exists a maximal wealth process $X'$ such that\n\\begin{equation}\\nonumbe\nX'_t \\geq |e_t|, \\quad for~every~t\\in[0,T],\\quad\\Pas.\n\\end{equation}\nMoreover, Assumption \\ref{asEnd} and \\eqref{NFLVR} imply that all the assertions of Lemma \\ref{lemma1} hold.\n\\end{Assumption}\n\n\\begin{Remark}\nIf $\\theta _0 = \\{T\\}$, then Assumption~\\ref{asEnd} is equivalent to the assumptions on endowment in~\\cite{HK04} (for the case of one-dimensional random endowment).\n\\end{Remark}\n\n\n\nThe preferences of an economic agent are modeled with a\n\\textit{utility stochastic field}\n$U=U(t, \\omega, x):[0,T]\\times\\Omega\\times[0,\\infty)\\to\\mathbb R\\cup \\{-\\infty\\}$.\nWe assume that $U$\nsatisfies the conditions below.\n\\begin{Assumption}\\label{assumptionOnU}\n \nFor every $(t, \\omega)\\in[0, T]\\times\\Omega$, the function $x\\to\n U(t, \\omega, x)$ is strictly concave, increasing, continuously\n differentiable on $(0,\\infty)$ and satisfies the Inada conditions:\n \\begin{equation}\\nonumbe\n \\lim\\limits_{x\\downarrow 0}U'(t, \\omega, x) =\n \\infty \\quad \\text{and} \\quad \\lim\\limits_{x\\to\n \\infty}U'(t, \\omega, x) = 0,\n \\end{equation}\n where $U'$ denotes the partial derivative with respect to the third argument.\n At $x=0$ we suppose, by continuity, $U(t, \\omega, 0) = \\lim\\limits_{x \\downarrow 0}U(t, \\omega, x)$, {which}\n may be $-\\infty$.\nFor every\n $x\\geq 0$ the stochastic process $U\\left( \\cdot, \\cdot, x \\right)$ is\n optional. Below, following the standard convention, we will not write $\\omega$ in $U$.\n\\end{Assumption}\n\n\n\nThe agent can control investment and consumption. {The goal is to maximize}\n expected utility. The value function $u$ is\ndefined as:\n\\begin{equation}\\label{primalProblem}\nu(x, q) {:=} \\sup\\limits_{c\\in\\mathcal A(x, q)}\\mathbb\nE\\left[\\int_0^T U(t, c_t)d\\kappa_t\\right], \\quad\n(x, q)\\in\\mathcal K.\n\\end{equation}\nIn \\eqref{primalProblem}, we use the convention\n\\begin{displaymath}\n \\mathbb{E}\\left[\n \\int_0^TU(t,c_t)d\\kappa_t \\right] {:=} -\\infty\n \\quad \\text{if} \\quad \\mathbb{E}\\left[ \\int_0^TU^{-}(t,c_t)d\\kappa_t\n \\right]= \\infty.\n\\end{displaymath}\nHere and below, $W^{-}$ and $W^{+}$ denote the negative and positive parts of a\nstochastic field $W$, respectively.\n\n\n We employ duality techniques to obtain the standard conclusions of the utility maximization theory. We first define the convex conjugate stochastic field\n\n\\begin{equation}\\label{defV}\nV(t, y) {:=} \\sup\\limits_{x>0}\\left(\nU(t, x)-xy\\right),\\quad (t, y)\\in[0,T]\\times [0,\\infty),\n\\end{equation}\n and then observe that \n$-V$ satisfies Assumption \\ref{assumptionOnU}.\nIn order to construct the feasible set of the dual problem, we define the set\n$\\mathcal L$ as the polar cone of $-\\mathcal K$:\n\\begin{equation}\\label{defL}\n\\mathcal L {:=} \\left\\{(y,r)\\in\\mathbb R\\times\\mathbb L^1(dK): xy + \\int_0^Tq_sr_sdK_s \\geq\n0{\\rm~for~every~}(x,q)\\in\\mathcal K\\right\\}.\n\\end{equation}\n\\begin{Remark}\nUnder the conditions \\eqref{finClock}, \\eqref{NFLVR} and Assumption \\eqref{asEnd}, the set \n$\\mathcal L$ is non-empty. By definition, it is closed in $\\mathbb R\\times\\mathbb L^1(d K)$-norm and $\\sigma(\\mathbb R\\times\\mathbb L^1(dK), \\mathbb R\\times\\mathbb L^{\\infty}(dK))$ topologies. Also, as shown later, the set $\\mathcal L = (-\\mathring{\\mathcal K})^o$, i.e. the polar of $-\\mathring{\\mathcal K}$.\n\\end{Remark}\nBy ${\\mathcal Z}$, we denote the set of c\\`adl\\`ag densities of\nequivalent local martingale measures:\n\\begin{equation}\\label{defZ}\n{\\mathcal Z} {:=} \\left\\{{\\rm c\\grave adl\\grave ag}~\\left(\\frac{d\\mathbb Q_t}{d\\mathbb P_t}\\right)_{t\\in[0,T]}:\\quad \\mathbb Q\\in\\mathcal M\\right\\}.\n\\end{equation}\nLet us denote by $\\mathbb L^0 = \\mathbb L^0(d\\kappa \\times \\mathbb P)$ the linear space of (equivalence classes of) real-valued optional processes on the stochastic basis $(\\Omega, \\mathcal F, (\\mathcal F_t)_{t\\in[0,T]}, \\mathbb P)$ which we equip with the topology of convergence in measure $(d\\kappa \\times \\mathbb P)$. \nFor each $y \\geq 0$ we define\n\\begin{equation}\\label{oldY}\n\\begin{array}{c}\n \\hspace{-15mm}\\mathcal Y(y) {:=} {\\cl}\\left\\{ Y: ~Y{\\rm~is~c\\grave{a}dl\\grave{a}g~adapted~and~}\\right. \\\\\n \\hspace{35mm}\\left.0\\leq Y\\leq yZ ~\\left(d\\kappa\\times\\mathbb\nP\\right){\\rm\n~a.e.~for~some~}Z\\in{\\mathcal Z}\\right\\},\\\\\n\\end{array}\n\\end{equation}\nwhere the closure is taken in\n$\\mathbb L^0$. \nNow we are ready to set the\ndomain of the dual problem:\n\\begin{equation}\\label{defY}\n\\begin{array}{rcl}\n\\mathcal Y(y,r)&{:=} &\\left\\{\n Y: Y \\in\\mathcal Y(y)\n~{\\rm and}\\right.\\\\\n&&\n\\quad\\mathbb E\\left[\\int_0^Tc_sY_sd\\kappa_s \\right]\\leq xy + \\int_0^Tq_sr_sdK_s,\\\\\n&&\\left.\\quad{\\rm for~every}~(x,q)\\in\\mathcal K {\\rm ~and~}c\\in\\mathcal A(x,q)\\right\\}\\\\\n\\end{array}\n\\end{equation}\nNote that the definition \\eqref{defY} requires that every element of $\\mathcal Y(y,r)$ is in $\\mathcal Y(y)$, $y\\geq 0$. Also, for every $(y,r)\\in\\mathcal L$, $\\mathcal Y(y,r)\\neq \\emptyset$, since $0\\in\\mathcal Y(y,r)$. \n\n We can now state the dual optimization problem:\n\\begin{equation}\\label{dualProblem}\nv(y, r) {:=} \\inf\\limits_{Y\\in\\mathcal Y(y, r)}\\mathbb\nE\\left[\\int_0^T V(t, Y_t)d\\kappa_t\\right],\\quad(y, r)\\in\\mathcal L,\n\\end{equation}\nwhere we use the convention:\n\\begin{displaymath}\n \\mathbb{E}\\left[ \\int_0^TV(t,Y_t\n )d\\kappa_t \\right] {:=} \\infty\n \\quad \\text{if} \\quad \\mathbb{E}\\left[ \\int_0^TV^{+}(t, Y_t )d\\kappa_t\n \\right] = \\infty.\n\\end{displaymath}\nAlso, we set \n\\begin{equation}\\label{3313}\nv(y,r) {:=} \\infty~for~(y,r)\\in\\mathbb R\\times \\mathbb L^1(dK)\\backslash\n \\mathcal L\\quad and \\quad u(x,q){:=} -\\infty~for~(x,q)\\in\\mathbb R\\times \\mathbb L^{\\infty}(dK)\\backslash\\mathcal K.\\\\ \n\\end{equation}\nWith this definition, it will be shown below in Theorem \\ref {mainTheorem} that $u<\\infty$ and $v>-\\infty$ everywhere, so $u$ and $v$ are proper functions in the language of convex analysis. \nLet us recall that in the absence of random endowment, the dual value function is defined as\n\\begin{equation}\\nonumbe\n\\tilde w(y) {:=} \\inf\\limits_{Y\\in\\mathcal Y(y)}\\mathbb E\\left[\\int_0^TV(t, Y_s)d\\kappa_s \\right],\\quad y>0,\n\\end{equation}\nwhereas the primal value function is given by\n\\begin{equation}\\label{defw}\nw(x) {:=} u(x,0),\\quad x>0.\n\\end{equation}\n\\section{Existence, uniqueness, and biconjugacy}\\label{secProofs}\n\\begin{Theorem}\\label{mainTheorem}\nLet (\\ref{finClock}) and (\\ref{NFLVR}), Assumptions~\\ref{asEnd}\nand \\ref{assumptionOnU} hold true and\n\\begin{equation}\\tag{finValue}\\label{finValue}\nw(x) > -\\infty\\quad for~every~x>0\\quad and \\quad \n \\widetilde w(y) <\\infty\\quad for~every~y>0.\n \\end{equation}\n Then we have:\n\n$(i)$ $u$ is finite-valued on $\\mathring{\\mathcal K}$ and $u<\\infty$ on $\\mathbb R\\times \\mathbb L^\\infty(dK)$. The dual value function $v$ satisfies $v>-\\infty$ on $\\mathbb R\\times\\mathbb L^1(dK)$, and the set $\\{v<\\infty\\}$ is a nonempty convex subset of $\\mathcal L$, whose closure in $\\mathbb R\\times\\mathbb L^1(dK)$ equals to $\\mathcal L$. \n\n$(ii)$ $u$ is concave, proper, and upper semicontinuous with respect to the norm-topology of $\\mathbb R\\times\\mathbb L^\\infty(dK)$ and the weak-star topology $\\sigma (\\mathbb R\\times\\mathbb L^\\infty(dK), \\mathbb R\\times\\mathbb L^1(dK))$. For every $(x,q)\\in \\{u>-\\infty\\}$, there exists a unique solution to (\\ref{primalProblem}). Likewise, $v$ is convex, proper, and lower semicontinuous with respect to the norm-topology of $\\mathbb R\\times \\mathbb L^1(dK)$ and the the weak topology $\\sigma(\\mathbb R\\times\\mathbb L^1(dK),\\mathbb R\\times\\mathbb L^{\\infty}(dK))$. For every $(y,r)\\in \\{v<\\infty\\}$, there exists a unique solution to (\\ref{dualProblem}). \n\n$(iii)$ The functions $u$ and $v$ satisfy the biconjugacy relations\n\\begin{equation}\\nonumbe\n\\begin{array}{rcll}\nu(x,q) &=& \\inf\\limits_{(y,r)\\in \\mathcal L}\\left(v(y,r) + xy + \\int_0^Tr_sq_sdK_s\\right),& (x,q)\\in \\mathcal K,\\\\\nv(y,r) &=& \\sup\\limits_{(x,q)\\in \\mathcal K}\\left(u(x,q) - xy - \\int_0^Tr_sq_sdK_s\\right),& (y,r)\\in \\mathcal L.\\\\\n\\end{array}\n\\end{equation}\n\\end{Theorem}\n\n\n\n\\begin{Lemma}\\label{lemma1}\nLet (\\ref{NFLVR}) and Assumption~\\ref{asEnd} hold. Then we have:\n\n$(i)$ for every $x>0$, $(x,0)$ belongs to $\\mathring{\\mathcal K}$ (in particular, $\\mathring{\\mathcal K}\\neq \\emptyset$),\n\n$(ii)$ for every $q\\in\\mathbb L^\\infty(dK)$, there exists $x>0$ such that $(x,q)\\in\\mathring{\\mathcal K}$,\n\n$(iii)$ $\\sup\\limits_{\\mathbb Q\\in\\mathcal M}\\mathbb {E^Q}\\left[\\int_0^T|e_s|d\\kappa_s \\right] <\\infty$,\n\n$(iv)$ there exists a nonnegative maximal wealth process $X^{''}$, such that\n\\begin{displaymath}\nX^{''}_T \\geq \\int_0^T|e_s|d\\kappa_s,\\quad \\Pas,\n\\end{displaymath}\n\n$(v)$ there exists a nonnegative maximal wealth process $X^{''}$, such that\n\\begin{equation}\\nonumbe\nX^{''}_t \\geq \\int_0^t|e_s|d\\kappa_\ns,\\quad t\\in[0,T],~\\mathbb P-a.s.\n\\end{equation}\n\n\\end{Lemma}\n\\begin{proof}\nFirst, via \\cite[Proposition I.4.49]{JS} and \\eqref{finClock}, we get\n\\begin{equation}\\nonumber\n\\begin{array}{c}\n\\int_0^T |e_s|d\\kappa_s \\leq \\int_0^T X'_sd\\kappa_s = -\\int_0^T\\kappa_{s-}dX'_s + \\kappa_TX'_T \\\\ \\leq -\\int_0^T\\kappa_{s-}dX'_s +AX'_T =AX'_0 + \\int_0^T(A - \\kappa_{s-})dX'_s.\\\\\n\\end{array}\n\\end{equation}\nTherefore, the exists a self-financing wealth process $\\bar X$, such that\n$$\\int_0^T |e_s|d\\kappa_s \\leq \\bar X_T.$$\nConsequently, \\cite[Theorem 2.3]{DS97} asserts the existence of a nonnegative {\\it maximal} process $X^{''}$, such that\n$$\\int_0^T |e_s|d\\kappa_s \\leq X_T^{''},$$\ni.e., $(iv)$ holds. Therefore $(iii)$ is valid as well, by \\cite[Theorem 5.12]{DS1998}. \n\nRelation $(iv)$ and \\eqref{NFLVR} imply $(v)$.\nTo prove $(i)$ and $(ii)$, without loss of generality, we will suppose that in $(v)$, $X^{''}_0> 0.$ Let $q\\in\\mathbb L^{\\infty}(dK)$ and $\\varepsilon >0$ be fixed. Let us define\n\\begin{equation}\\nonumbe\nx(q){:=} \\left(||q||_{\\mathbb L^{\\infty}(dK)} + 2\\varepsilon\\right)X^{''}_0.\n\\end{equation}\nWe claim that $(x(q), q)\\in\\mathring{\\mathcal K}.$\nLet \nus consider arbitrary \n\\begin{equation}\\label{454}\n|x'|\\leq \\varepsilon\\quad {\\rm and} \\quad q'\\in\\mathbb L^\\infty(dK)\n:~||q'||_{\\mathbb L^{\\infty}(dK)}\\leq {\\varepsilon},\n\\end{equation} \nand set\n\\begin{equation}\\label{455}\n\\tilde X_t {:=} \\left(||q||_{\\mathbb L^{\\infty}(dK)} + 2\\varepsilon+x'\\right)X^{''}_t,\\quad t\\in[0,T].\n\\end{equation} Then, by item $(v)$, we have\n\\begin{displaymath}\n\\begin{array}{c}\n\\tilde X_t\\geq \\left(||q||_{\\mathbb L^{\\infty}(dK)}+{\\varepsilon}\\right)X^{''}_t\\geq \\left(||q||_{\\mathbb L^{\\infty}(dK)}+{\\varepsilon}\\right)\\int_0^t|e_s|d\\kappa_s\\geq -\\int_0^t(q_s + q'_s)|e_s|d\\kappa_s.\n\\end{array}\n\\end{displaymath}\nWe deduce that $\\tilde X\\in\\mathcal X(\\left(||q||_{\\mathbb L^{\\infty}(dK)} + 2\\varepsilon+x'\\right)X^{''}_0, q+q')$. In particular, $$\\mathcal X(\\left(||q||_{\\mathbb L^{\\infty}(dK)} + 2\\varepsilon+x'\\right)X^{''}_0, q+q')\\neq \\emptyset.$$ As $x'$ and $q'$ are arbitrary elements satisfying \\eqref{454}, we deduce that $(x(q), q)\\in\\mathring{\\mathcal K}$. This proves $(ii)$. \n\n\nIn order to show $(i)$, first we observe that $\\mathring{\\mathcal K}$ is a convex cone. Therefore, it suffices to prove that, for a given $\\varepsilon>0$, we have\n\\begin{equation}\\label{2121}\n\\left(2\\varepsilon X^{''}_0, 0\\right)\\in\\mathring{\\mathcal K}.\n\\end{equation} \nAgain, let us consider $x'$ and $q'$ satisfying \\eqref{454} and $\\tilde X$ satisfying \\eqref{455} for $q\\equiv 0$. \nThen for every $t\\in[0,T]$, we have\n$$\\tilde X_t\\geq {\\varepsilon}X^{''}_t\\geq {\\varepsilon}\\int_0^t|e_s|d\\kappa_s\\geq -\\int_0^tq'_s|e_s|d\\kappa_s.$$\nThus, $\\mathcal X(\\left(2\\varepsilon+x'\\right)X^{''}_0, q')\\neq \\emptyset$. Consequently, as $x'$ and $q'$ are arbitrary elements satisfying \\eqref{454}, \\eqref{2121} holds and so is $(i)$. \n This completes the proof of the lemma. \n\\end{proof}\n\\begin{Remark}\nA close look at the proofs shows that the conclusions of Theorem \\ref{mainTheorem} also hold if instead of Assumption \\ref{asEnd}, we impose any of the equivalent assertions $(iii) - (v)$ of Lemma \\ref{lemma1}.\n\\end{Remark}\n\n\n\\subsection{Characterization of the primal and dual domains}\\label{subSecCharDualDomain}\nThe polar, $A^{o}$, of a nonempty subset $A$ of $\\mathbb R\\times\\mathbb L^\\infty(dK)$, is\n the subset of\n $\\mathbb R\\times\\mathbb L^1(dK)$, defined by\n$$A^o{:=} \\left\\{(y,r)\\in \\mathbb R\\times \\mathbb L^1(dK):~xy + \\int_0^T q_sr_sdK_s \\leq 1, \\quad for~every~(x,q)\\in A\\right\\}.$$\nThe polar of a subset of $\\mathbb R\\times\\mathbb L^1(dK)$ is defined similarly.\n\\begin{Proposition}\\label{prop1} Under Assumption \\ref{asEnd} and conditions \\eqref{finClock} and \\eqref{NFLVR}, we have:\n\n$(i)$ Let $(x,q)\\in\\mathbb R\\times\\mathbb L^{\\infty}(dK)$. Then $c\\in \\mathcal A(x,q)$ (thus, $\\mathcal A(x,q)\\neq \\emptyset$ so $(x,q)\\in \\mathcal K$ ) if and only if\n\\begin{equation}\\nonumber\n\\mathbb E\\left[\\int_0^T c_sY_sd\\kappa_s \\right]\\leq xy + \\int_0^Tq_sr_sdK_s,\n~~ {\\rm for~every}~(y,r)\\in \\mathcal L{\\rm ~ and }~ Y\\in\\mathcal Y(y,r).\n\\end{equation}\n\n$(ii)$ Likewise, for $(y,r)\\in \\mathbb{R}\\times \\mathbb{L}^1 (dK)$ we have $(y,r)\\in \\mathcal L$ and $Y\\in \\mathcal Y(y,r)$ if and only if\n\\begin{equation}\\nonumber\n\\mathbb E\\left[\\int_0^T c_sY_sd\\kappa_s \\right]\\leq xy + \\int_0^Tq_sr_sdK_s,\n~~ {\\rm for~every}~(x,q)\\in\\mathcal K{\\rm ~ and }~c\\in\\mathcal A(x,q).\n\\end{equation}\n\nWe also have $\\mathcal K = (-\\mathcal L)^{o} = \\cl{\\mathring{\\mathcal K}}$, where the closure is taken both in norm $\\mathbb R\\times \\mathbb L^1(dK)$ and $\\sigma(\\mathbb R\\times\\mathbb L^\\infty(dK), \\mathbb R\\times\\mathbb L^1(dK))$ topologies.\n\\end{Proposition}\n\\begin{Remark}\\label{remNonEmpty}\nIt follows from Proposition \\ref{prop1} that \nfor every $(x,q)\\in\\mathcal K, \\mathcal A(x,q)\\neq \\emptyset$ as $0\\in\\mathcal A(x,q)$. Likewise, \nfor every $(y,r)\\in \\mathcal L, \\mathcal Y(y,r)\\neq \\emptyset$ as $0\\in\\mathcal Y(y,r)$. Moreover, for every $(x,q)\\in \\mathring{\\mathcal K}$, each of the sets $\\mathcal A(x,q)$ and $\\bigcup\\limits_{(y,r)\\in \\mathcal L: ~xy + \\int_0^Tr_sq_sdK_s \\leq 1}\\mathcal Y(y,r)$ contain a strictly positive element, see Lemma \\ref{6-12-1} below. \n\\end{Remark}\nThe proof of Proposition~\\ref{prop1} will be given via several lemmas.\nLet \n\\begin{equation}\\nonumbe\n\\begin{array}{l}\n\\mathcal M'~be~the~set~of~equivalent~local~martingale~measures,~under~which\\\\X^{''}~(from~Lemma~\\ref{lemma1},~item~(v))~is~a~uniformly~integrable~martingale. \n\\end{array}\n\\end{equation}\nNote that by \\cite[Theorem 5.2]{DS97}, $\\mathcal M'$ is a nonempty, convex subset of $\\mathcal M$, which is also dense in $\\mathcal M$ in the total variation norm. \n\n\\begin{Remark}\nEven though the results in \\cite{DS97} are obtained under the condition that $S$ is a locally bounded process, they also hold without local boundedness assumption, see the discussion in \\cite[Remark 3.4]{HKS05}.\n\\end{Remark}\n\nLet $Z'$ denote the set of the corresponding c\\`adl\\`ag densities,~i.e., \n\\begin{equation}\\label{defZ'}\n\\begin{array}{rcl}\n\\mathcal Z' &{:=}& \\left\\{{\\rm c\\grave adl\\grave ag}~Z: Z_t=\\mathbb E\\left[ \\frac{d\\mathbb Q}{d\\mathbb P}|\\mathcal F_t\\right],~t\\in[0,T],~\\mathbb Q\\in\\mathcal M'\\right\\}.\\\\\n\\end{array}\n\\end{equation}\n We also set\n\\begin{equation}\\label{defUpsilon}\n\\Upsilon {:=} \\left\\{1_{[0,\\tau]}(t),t\\in[0,T]:~~\\tau\\in\\Theta\\right\\}.\n\\end{equation}\n\n\\begin{Lemma}\\label{6-16-1}\nLet the conditions of Proposition \\ref{prop1} hold, $\\mathbb Q\\in\\mathcal M'$, ${Z} = {Z}^{\\mathbb Q}$ be the corresponding element of $\\mathcal Z'$, and $\\Lambda \\in \\Upsilon$. Then there exists $r\\in \\mathbb L^1(dK)$, uniquely defined~by\n\\begin{equation}\\label{6-25-1}\n\\int_0^tr_sdK_s = \\mathbb E\\left[\\int_0^t\\Lambda_s{Z}_se_sd\\kappa_s \\right],\\quad t\\in[0,T],\n\\end{equation}\nsuch that $(1,r)\\in \\mathcal L$ and ${Z}\\Lambda = ({Z}_t\\Lambda_t)_{t\\in[0,T]}\\in\\mathcal Y(1,r)$.\n\\end{Lemma}\n\\begin{proof}\nLet $(x,q)\\in\\mathcal K$ and $c\\in\\mathcal A(x,q)$.\nThen there exists \n$X\\in\\mathcal X(x,q)$, such that\n\\begin{equation}\\label{246}\n\\int_0^\\tau c_sd\\kappa_s\\leq X_\\tau + \\int_0^\\tau q_se_sd\\kappa_s,\\quad for~every~\\tau\\in\\Theta.\n\\end{equation}\nIn particular, \\eqref{246} holds for the particular $\\tau$, such that $\\Lambda_t = 1_{[0,\\tau]}(t)$, $t\\in[0,T]$. \nLet $G_t{:=} \\int_0^t(q_se_s)^{+}d\\kappa_s$, $t\\in[0,T]$, where $(\\cdot)^{+}$ denotes the positive part. By \\cite[Theorem III.29, p. 128]{Pr}, $\\int_0^tG_{s-}d{Z}_s$, $t\\in[0,T]$, is a local martingale, so let $(\\sigma_n)_{n\\in\\mathbb N}$ be its localizing sequence. Then by the monotone convergence theorem, integration by parts formula, and the optional sampling theorem, we get\n\\begin{equation}\\label{281}\n\\begin{array}{c}\n\\mathbb E\\left[\\int_0^T {Z}_s(q_se_s)^{+}\\Lambda_sd\\kappa_s \\right] = \\lim\\limits_{n\\to\\infty}\\mathbb E\\left[\\int_0^{\\sigma_n\\wedge \\tau} {Z}_sdG \\right] \\\\\n= \\lim\\limits_{n\\to\\infty} \\mathbb E\\left[{Z}_{\\sigma_n\\wedge \\tau}G_{\\sigma_n\\wedge \\tau}\\right] \n-\\lim\\limits_{n\\to\\infty} \\mathbb E\\left[ \\int_0^{\\sigma_n\\wedge \\tau}G_{s-}d{Z}_s\\right] \n\\\\=\n\\lim\\limits_{n\\to\\infty} \\mathbb E\\left[{Z}_{\\sigma_n\\wedge \\tau}G_{\\sigma_n\\wedge \\tau} \\right] = \\lim\\limits_{n\\to\\infty} \\mathbb E^{\\mathbb Q}\\left[G_{\\sigma_n\\wedge \\tau} \\right] \n = \\mathbb E^{\\mathbb Q}\\left[G_\\tau \\right],\\\\\n\\end{array}\n\\end{equation}\nwhere in the last equality we used the monotone convergence theorem again.\nHere finiteness of $\\mathbb E^{\\mathbb Q}\\left[G_\\tau \\right]=\\mathbb E^{\\mathbb Q}\\left[\\int_0^\\tau(q_se_s)^{+}d\\kappa_s \\right]$ follows from Assumption \\ref{asEnd} via Lemma \\ref{lemma1}, part (iii). \nIn a similar manner, we can show that\n\\begin{equation}\\label{283}\n\\begin{array}{rcl}\n\\mathbb E\\left[\\int_0^T(q_se_s)^{-}{Z}_s\\Lambda_sd\\kappa_s\\right] &= &\\mathbb E^{\\mathbb Q}\\left[ \\int_0^\\tau (q_se_s)^{-}d\\kappa_s\\right]<\\infty,\\\\\n\\mathbb E\\left[\\int_0^Tc_s{Z}_s\\Lambda_sd\\kappa_s\\right] &=& \\mathbb E^{\\mathbb Q}\\left[ \\int_0^\\tau c_sd\\kappa_s\\right]<\\infty.\\\\\n\\end{array}\n\\end{equation}\n\\cite[Lemma 4]{HK04} applies here and asserts that $X$ in \\eqref{246} is a supermartingale under $\\mathbb Q$. \nTherefore, from \\eqref{246}, using \\eqref{281} and \\eqref{283}, and taking expectation under $\\mathbb Q$, we get\n\\begin{equation}\\label{6-22-2}\n \\mathbb E\\left[ \\int_0^Tc_s\\Lambda_s{Z}_sd\\kappa_s\\right]\\leq x + \\mathbb E\\left[\\int_0^Tq_se_s{Z}_s\\Lambda_sd\\kappa_s\\right].\n\\end{equation}\nLet us define\n\\begin{equation}\\nonumbe\nR_t {:=} \\mathbb E\\left[\\int_0^t\\Lambda_s{Z}_se_sd\\kappa_s \\right],\\quad t\\in[0,T].\n\\end{equation}\nUsing the monotone class theorem, we obtain\n\\begin{equation}\\label{6-25-2}\n\\mathbb E\\left[\\int_0^T\\tilde q_se_s{Z}_s\\Lambda_sd\\kappa_s\\right]\n=\n\\int_0^T\\tilde q_sdR_s,\\quad for~every\\quad\\tilde q\\in\\mathbb L^\\infty(dK).\n\\end{equation}\nWe claim that\n$dR$ is absolutely continuous with respect to $dK$. First, using the $\\pi-\\lambda$ theorem, one can show that for every Borel-measurable subset $A$ of $[0,T]$, we have\n$$K(A) = \\mathbb E\\left[\\int_0^T 1_A(t)d\\kappa_t \\right]\\quad and \\quad R(A) = \\mathbb E\\left[\\int_0^T1_A(t)\\Lambda_t{Z}_te_td\\kappa_t\\right].$$\nThus, if for some $A$, $K(A) = 0$, then $\\int_0^T 1_A(t)d\\kappa_t=0$ a.s. and $\\kappa^A_t {:=} \\int_0^t 1_A(s)d\\kappa_s$, $t\\in[0,T]$, satisfies $\\kappa^A_T = 0$ a.s. and $\\int_0^T \\Lambda_t{Z}_te_t d\\kappa^A_t = \\int_0^T\\Lambda_t{Z}_te_t1_A(t)d\\kappa_t = 0$ a.s. \n\n\nAs $dR$ is absolutely continuous with respect to $dK$, \nthere exists a unique $r\\in\\mathbb L^1(dK)$, such that \\eqref{6-25-1}\nholds.\nSince the left-hand side in (\\ref{6-22-2}) is nonnegative and since $(x,q)$ is an arbitrary element of $\\mathcal K$, we deduce from the definition of $\\mathcal L$, \\eqref{defL}, that $(1,r)\\in \\mathcal L$. Finally, it follows from \\eqref{6-22-2} and \\eqref{6-25-2}\nthat ${Z}\\Lambda \\in\\mathcal Y(1,r)$. This completes the proof of the~lemma.\n\\end{proof}\n\\begin{Remark}\\label{Remark-Mihai-2}\nThe natural convexification of the set $\\Upsilon$ consists of non-negative left-continuous decreasing and adapted processes $D$ such that $D_{\\theta _0}=1$. In the context of utility-maximization constraints (and for $\\theta _0=0$), this is follows from \\cite{HePages} and \\cite{ElKarouiJeanblanc}. \nWe investigate convexification of the constraints in the later Sections \\ref{secStructureOfDualDomain} and \\ref{secSlackness}, where will extend Lemma \\ref{6-16-1} to a more general set of decreasing processes than $\\Lambda$ that drives the dual domain and that allows for the multiplicative decomposition of the dual minimizer\n\\end{Remark}\n\\begin{Corollary}\\label{cor1}\nLet the conditions of Proposition \\ref{prop1} hold, $\\mathbb Q\\in\\mathcal M'$, ${Z}$ be the c\\'adl\\'ag modification of the density process $\\mathbb E\\left[ \\frac{d\\mathbb Q}{d\\mathbb P}|\\mathcal F_t\\right]$, $t\\in[0,T]$. Then there exists $r\\in\\mathbb L^1(dK)$, such that\n ${Z}\\in\\mathcal Y(1,r)$,\nwhere\n$$\n\\int_0^tr_sdK_s = \\mathbb E\\left[\\int_0^t{Z}_se_sd\\kappa_s \\right],\\quad t\\in[0,T],\n$$\nand $(1,r)\\in \\mathcal L$.\n\\end{Corollary}\n\\begin{proof}[Proof of Lemma \\ref{3311}] Let us fix $\\mathbb Q\\in\\mathcal M'$ and let $Z\\in\\mathcal Z'$ be the corresponding density process. As in Lemma \\ref{6-16-1}, we can show that there exists $r\\in\\mathbb L^1(dK)$, such that for every bounded and Borel measurable function $q$ on $[0,T]$ we have\n\\begin{equation}\\label{5192}\n\\int_0^tq_s r_sdK_s = \\mathbb E\\left[\\int_0^t q_sZ_s|e_s|d\\kappa_s \\right],\\quad t\\in[0,T].\n\\end{equation}\nLet $\\bar q {:=} |q^1 - q^2|$, then as $\\bar q = 0$, $dK$-a.e., we get\n\\begin{displaymath}\n\\mathbb E\\left[\\int_0^T \\bar q_s|e_s|Z_sd\\kappa_s \\right] = \\int_0^T \\bar q_sr_sdK = 0.\n\\end{displaymath}\nTherefore, using integration by parts and via \\eqref{5192}, we obtain\n$$\n0 = \\int_0^T \\bar q_sr_sdK = \\mathbb E\\left[\\int_0^T q_sZ_s|e_s|d\\kappa_s \\right] = \\mathbb E^{\\mathbb Q}\\left[\\int_0^T q_s|e_s|d\\kappa_s \\right].\n$$\nConsequently, $\\int_0^T q_s|e_s|d\\kappa_s = 0$, $\\mathbb Q$-a.s., and by the equivalence of $\\mathbb Q$ and $\\mathbb P$, also $\\mathbb P$-a.s. \nAs, by construction $\\int_0^t q_s|e_s|d\\kappa_s = 0$, $t\\in[0,T]$, is a nonnegative and non-decreasing process, whose terminal value is $0$, $\\mathbb P$-a.s., we conclude that it is indistinguishable from the $0$-valued process. The assertions of the lemma follows.\n\n\\end{proof}\n\\begin{Lemma}\\label{9-5-1}\nLet the conditions of Proposition \\ref{prop1} hold, $(x,q)\\in\\mathcal K$, and $c$ is a nonnegative optional process. Then $c\\in\\mathcal A(x,q)$ if and only if\n\\begin{equation}\\label{241}\n\\mathbb E\\left[\\int_0^Tc_s\\Lambda_s{Z}_sd\\kappa_s \\right]\\leq\nx + \\mathbb E\\left[\\int_0^Tq_se_s\\Lambda_s{Z}_sd\\kappa_s \\right] = x + \\int_0^Tq_sr_sdK_s,\n\\end{equation}\nfor every ${Z}\\in\\mathcal Z'$ and $\\Lambda \\in \\Upsilon$, where $r$ is given by (\\ref{6-25-1}).\n\\end{Lemma}\n\\begin{proof}\nLet $c\\in\\mathcal A(x,q)$. Then for every ${Z}\\in\\mathcal Z'$ and $\\Lambda \\in \\Upsilon$, the validity of \\eqref{241} follows from the definition of $\\mathcal A(x,q)$, integration by parts formula and supermartingale property of every $X\\in\\mathcal X(x,q)$ under every $\\mathbb Q\\in\\mathcal M'$, which in turn follows from \\cite[Lemma 4]{HK04}. \n\nConversely, let \\eqref{241} holds for every ${Z}\\in\\mathcal Z'$ and $\\Lambda \\in \\Upsilon$. \nThen, we have\n\\begin{equation}\\label{243}\n\\mathbb E\\left[\\int_0^T(c_s-q_se_s){Z}_s\\Lambda_sd\\kappa_s \\right]\\leq x,\n\\end{equation}\nwhich, in view of the definition of $\\Upsilon$ in \\eqref{defUpsilon}, localization, and integration by parts, implies that\n\\begin{displaymath}\n\\sup\\limits_{\\mathbb Q\\in\\mathcal M', \\tau\\in \\Theta}\\mathbb {E^Q}\\left[\\int_0^{\\tau}(c_s-q_se_s)d\\kappa_s \\right]\\leq x,\n\\end{displaymath}\nFor $X^{''}$ given by Lemma \\ref{lemma1}, item $(v)$, let us denote\n\\begin{equation}\\label{244}\nf_t{:=} ||q||_{\\mathbb L^\\infty(dK)} X^{''}_t +\\int_0^{t}(c_s-q_se_s)d\\kappa_s,\\quad t\\in[0,T].\n\\end{equation}\nIt follows from Assumption \\ref{asEnd} and item $(v)$ of Lemma \\ref{lemma1} \n\n that \n $f$ is a nonnegative process.\nWe observe that the proof of \\cite[Proposition 4.3]{K96} goes through, if we only take stopping times in $\\Theta$ and measures in $\\mathcal M'$. This proposition allows to conclude that there exists a nonnegative c\\`adl\\`ag process $V$, such that \n\\begin{equation}\\label{defV2}\nV_t = \\esssup\\limits_{\\tau\\in \\Theta: \\tau\\geq t,\\mathbb Q \\in\\mathcal M'}\\mathbb {E^Q}\\left[f_{\\tau}|\\mathcal F_t \\right],\\quad t\\in[0,T],\n\\end{equation}\n which is a supermartingale for every $\\mathbb Q \\in\\mathcal M'$. Therefore, by the density of $\\mathcal M'$ in $\\mathcal M$ in the norm topology of $\\mathbb L^1(\\mathbb P)$ and Fatou's lemma, $V$ is a supermartingale under every $\\mathbb Q\\in\\mathcal M$. Moreover, $V_0$ satisfies \n$$V_0\\leq x + ||q||_{\\mathbb L^\\infty(dK)} X^{''}_0,$$ by \\eqref{defV2}, \\eqref{243}, and by following the argument in the proof of Lemma \\ref{6-16-1}. \n\nWe would like to apply \nthe optional decomposition theorem of F\\\"olmer and Kramkov, \\cite[Theorem 3.1]{FK}. For this, we need to show that $V$ is a local supermartingale under every $\\mathbb Q$, such that every $X\\in\\mathcal X(1)$ is a $\\mathbb Q$-local supermartingale. However, $\\mathcal M$ is dense in the set of such measures in the norm topology of $\\mathbb L^1(\\mathbb P)$, by the results of Delbaen and Schachermayer, see \\cite[Proposition 4.7]{DS1998}. Therefore, the supermartingale property of $V$ under every such $\\mathbb Q$ follows from Fatou's lemma and supermartingale property of $V$ under every $\\mathbb Q\\in\\mathcal M$ established above. Therefore, by \\cite[Theorem 3.1]{FK}, we get\n$$V_t = V_0 + H\\cdot S_t - A_t,\\quad t\\in[0,T],$$\nwhere $A$ is a nonnegative increasing process that starts at $0$. \nSubtracting the constant $||q||_{\\mathbb L^\\infty(dK)} X^{''}_0$ from both sides of \\eqref{244}, we get\n\\begin{equation}\\nonumber\\begin{array}{c}\n\\int_0^{\\tau}(c_s-q_se_s)d\\kappa_s = f_\\tau -||q||_{\\mathbb L^\\infty(dK)} X^{''}_0 \\leq V_\\tau -||q||_{\\mathbb L^\\infty(dK)} X^{''}_0\\\\= V_0 -||q||_{\\mathbb L^\\infty(dK)} X^{''}_0 + H\\cdot S_\\tau-A_\\tau \\leq x + H\\cdot S_\\tau, \\quad \\tau\\in\\Theta,\n\\end{array}\\end{equation}\nwhere $x + H\\cdot S$ is acceptable, since $V$ is nonnegative. \nConsequently, \n$c\\in\\mathcal A(x,q)$. This completes the proof of the lemma.\n\\end{proof}\n\n\n\\begin{proof}[Proof of Proposition~\\ref{prop1}]\nThe assertions of item $(i)$ follow from Lemma~\\ref{9-5-1}. \nIt remains to show that the affirmations of item $(ii)$ hold.\nFix a $(y,r)\\in \\mathcal L$. If $Y\\in\\mathcal Y(y,r)$, $(ii)$ follows from the definition of $\\mathcal Y(y,r)$. Conversely,\nif $(ii)$ holds for a nonnegative process $Y$, then since $(x,0)\\in \\mathcal K$ for every $x>0$, we have \n$$\\mathbb E\\left[\\int_0^Tc_sY_sd\\kappa_s \\right]\\leq 1\\quad {\\rm for~every~}c\\in\\mathcal A\\left(\\frac{1}{y},0\\right).$$\nVia \\cite[Proposition 4.4]{MostovyiNec}, we deduce that $Y\\in\\mathcal Y(y)$ and is such that $(ii)$ holds. Therefore,\n$Y\\in\\mathcal Y(y,r)$. \n\nWe have $\\mathring{\\mathcal K}\\neq \\emptyset$, where the interior is taken with respect to the norm-topology. According to Proposition \\ref{prop1} (i) we have also $\\mathcal K = (-\\mathcal L)^{o}$. The set $\\mathcal{K}$, as the polar of $\\mathcal L$, is convex and closed both in (strong) $\\mathbb R\\times \\mathbb L^{\\infty}(dK)$ and $\\sigma(\\mathbb R\\times\\mathbb L^\\infty(dK), \\mathbb R\\times\\mathbb L^1(dK))$ topologies. Having non-empty strong interior, we obtain $\\mathcal K = \\cl{\\mathring{\\mathcal K}}$ where the closure is in the strong-topology. \nSince $\\mathcal{K}$ is also closed in the weaker $\\sigma(\\mathbb R\\times\\mathbb L^\\infty(dK), \\mathbb R\\times\\mathbb L^1(dK))$ topology we obtain that \n$$\\mathcal K = (-\\mathcal L)^{o} = \\cl{\\mathring{\\mathcal K}},$$ where the closure is taken in both\ntopologies.\n\\end{proof}\n\n\\subsection{Preliminary properties of the value functions, in particular finiteness. }\n\nLet $\\mathbf L^0_{+}$ be the positive orthant of $\\mathbf L^0$. The polar of a set $A \\subseteq \\mathbf L^0_{+}$ is defined as\n\\begin{displaymath} \nA^o{:=} \\left\\{ c\\in\\mathbf L^0_{+}:~\\mathbb E\\left[\\int_0^T c_sY_s d\\kappa_s\\right]\\leq 1,\\quad for~every~Y \\in A\\right\\}.\n\\end{displaymath}\n\nWe recall that the sets $\\mathcal Z$ and $\\mathcal Z'$ are defined in \\eqref{defZ} and \\eqref{defZ'}, respectively. \n \n\n\\begin{Lemma}\\label{finitenessOverZ'}\nUnder the conditions of Theorem~\\ref{mainTheorem}, we have\n\\begin{equation}\\label{bipolarZ'}\n(\\mathcal Z')^{oo} = \\mathcal Y(1)\n\\end{equation}\nand\n\\begin{equation}\\label{2145}\n\\tilde w(y) = \\inf\\limits_{Y\\in\\mathcal Z'}\\mathbb E\\left[\\int_0^TV(t, yY_s)d\\kappa_s \\right]<\\infty, \\quad y>0.\n\\end{equation}\n\\end{Lemma}\n\\begin{proof}\nIt follows from\n \\cite[Theorem 5.2]{DS97}\n that $\\mathcal Z'$ is dense in $\\mathcal Z$ \n in $\\mathbf L^0$. \nIt follows from Fatou's lemma that $$(\\mathcal Z')^{o} = \\mathcal Z^{o}.$$\nTherefore, \\cite[Lemma 4.2 and Proposition 4.4]{MostovyiNec} imply \\eqref{bipolarZ'}. \n\nOne can show that $\\mathcal Z'$ is closed under countable convex combinations, where the martingale property follows from the monotone convergence theorem \\and the c\\`adl\\`ag structure of the limit is guaranteed by \\cite[Theorem VI.18]{DelMey82}, see also \\cite[ Proposition 5.1]{KLPO14} for more details in similar settings. \nNow, \\eqref{2145} follows (up to a notational change) from \\cite[Theorem 3.3]{MostovyiNec}. This completes the proof of the lemma.\n\\end{proof}\n\\begin{Lemma}\\label{6-17-1} Under the conditions of Theorem~\\ref{mainTheorem}, \nfor every $(x,q)\\in\\mathcal K$ and $(y,r)\\in \\mathcal L$, we have\n$$u(x,q) \\leq v(y,r) + xy +\\int_0^Tr_sq_sdK_s.$$\n\\end{Lemma} \n\\begin{proof}\nFix an arbitrary $(x,q)\\in\\cl\\mathcal K$, $c\\in\\mathcal A(x,q)$ as well as $(y,r) \\in \\mathcal L$, $Y\\in\\mathcal Y(y,r)$. Using Proposition\n\\ref{prop1} and \\eqref{defV}, we get\n\\begin{displaymath}\n\\begin{array}{rcl}\n\\mathbb E\\left[\\int_0^T U (t, c_s)d\\kappa_s \\right] &\\leq &\\mathbb E\\left[\\int_0^T U(t, c_s)d\\kappa_s \\right] + xy + \\int_0^Tr_sq_sdK_s - \\mathbb E\\left[\\int_0^T c_sY_sd\\kappa_s \\right]\\\\\n&\\leq &\\mathbb E\\left[\\int_0^T V(t, Y_s)d\\kappa_s \\right] + xy + \\int_0^Tr_sq_sdK_s. \\\\\n\\end{array}\n\\end{displaymath}\nThis implies the assertion of the lemma.\n\\end{proof}\nFor every $(x,q)$ in $\\mathcal K$, we define\n\\begin{equation}\\label{auxiliarySets}\n\\begin{array}{rcl}\n\\mathcal B(x,q)&{:=}& \\left\\{(y,r)\\in \\mathcal L: ~xy + \\int_0^Tr_sq_sdK_s \\leq 1 \\right\\},\\\\\n\\mathcal D(x,q)&{:=}& \\bigcup\\limits_{(y,r)\\in \\mathcal B(x,q)}\\mathcal Y(y,r).\\\\\n\\end{array}\n\\end{equation}\nThe subsequent lemma established boundedness of $\\mathcal B(x,q)$ for $(x,q)$ in $\\mathring{\\mathcal K}$ in $\\mathbb R\\times \\mathbb L^1(dK)$.\n\\begin{Lemma}\n\\label{7-17-1}\nUnder the conditions of Theorem~\\ref{mainTheorem}, for every $(x,q)\\in\\mathring{\\mathcal K}$, \n$\\mathcal B(x,q)$ is bounded\n in $\\mathbb R\\times \\mathbb L^1(dK)$.\n\\end{Lemma}\n\\begin{proof}\nFix an $(x,q)\\in\\mathring{\\mathcal K}$. Then there exists $\\varepsilon > 0$, such that\nfor every\n\\begin{equation}\\label{7-20-1}\n|x'|\\leq\\varepsilon\\quad and \\quad||q'||_{\\mathbb L^{\\infty}} \\leq\\varepsilon,\n\\end{equation}\n we have\n$(x + x', q+q')\\in{\\mathring{\\mathcal K}}.$\nLet us fix an arbitrary $(y,r)\\in \\mathcal B(x,q)$. Then for every $(x',q')$ satisfying (\\ref{7-20-1}), by the definitions of $\\mathcal L$ and $\\mathcal B(x,q)$, respectively,~we~get \n\\begin{displaymath}\n\\begin{array}{rcl}\nxy + \\int_0^Tq_sr_sdK_s + x'y + \\int_0^Tq'_sr_sdK_s &\\geq& 0, \\\\\nxy + \\int_0^Tq_sr_sdK_s &\\leq& 1, \\\\\n\\end{array}\n\\end{displaymath}\nwhich implies that\n\\begin{equation}\\label{7-17-2}\n-x'y - \\int_0^Tq'_sr_sdK_s \\leq xy + \\int_0^Tq_sr_sdK_s \\leq 1.\n\\end{equation}\nTaking \n\\begin{displaymath}\nx' = -\\varepsilon\\quad{and}\\quad q'\\equiv 0,\n\\end{displaymath}\nwe deduce from (\\ref{7-17-2}) that $y\\leq \\frac{1}{\\varepsilon}$. Also, by the definition of $\\mathcal L$ and Lemma \\ref{lemma1}, item $(i)$, $y\\geq 0$.\nIn turn, setting\n\\begin{displaymath}\nx' = 0\\quad{and}\\quad q' = -\\varepsilon 1_{\\{r\\geq 0\\}} + \\varepsilon 1_{\\{r<0\\}},\n\\end{displaymath}\nwe obtain from (\\ref{7-17-2}) that $||r||_{\\mathbb L^1}\\leq \\frac{1}{\\varepsilon}$. This completes the proof of the lemma.\n\\end{proof}\n\\begin{Lemma}\\label{6-12-1}\nLet the conditions of Theorem~\\ref{mainTheorem} hold and $(x,q)$ be an arbitrary element of $\\mathring{\\mathcal K}$.\nThen, we have:\n\\newline\n$(i)$ $\\mathcal A(x,q)$ contains a strictly positive process.\n\\newline\n$(ii)$ The constant $\\bar y(x,q)$ given by \n\\begin{equation}\\label{defBary}\n\\bar y(x,q) {:=} \\frac{1}{|x| + ||q||_{\\mathbb L^\\infty(dK)}\\sup\\limits_{\\mathbb Q\\in\\mathcal M}\\mathbb {E^Q}\\left[\\int_0^T|e_s|d\\kappa_s \\right]},\n\\end{equation}\ntakes values in $(0,\\infty)$ and satisfies\n\\begin{equation}\\label{6-18-1}\n\\bar y(x,q)\\mathcal Z'\\subseteq \\mathcal D(x,q).\n\\end{equation}\nIn particular, for every $z>0$, $z\\mathcal D(x,q)$ contains a strictly positive process $Y$ such that \n\\begin{equation}\\label{6-16-2}\n\\mathbb E\\left[\\int_0^TV(s, Y_s)d\\kappa_s \\right] <\\infty.\n\\end{equation}\n\\end{Lemma}\n\\begin{proof}\nIn order to show $(i)$, we observe that the existence of a positive process in $\\mathcal A(x,q)$ follows from the fact that $(x-\\delta, q)\\in\\mathcal K$ for a sufficiently small $\\delta.$\nNow the constant-valued consumption $\\delta\/A>0$, where $A$ is the constant that dominates the terminal value of the stochastic clock $\\kappa$ in \\eqref{finClock}, is in $\\mathcal A(x,q)$.\n\nIn order to prove $(ii)$, \nlet us consider $\\bar y(x,q)$ given by \\eqref{defBary}.\nIt follows from Lemma \\ref{lemma1}, item $(iii)$, that $\\bar y (x,q)\\in (0,\\infty)$. \nFor this $\\bar y(x,q)$, using\nCorollary \\ref{cor1}, one can show~(\\ref{6-18-1}). This and Lemma~\\ref{finitenessOverZ'} (note that finiteness of $\\tilde w$ follows directly from \\eqref{finValue})\nimply that for every $z>0$, there exists a positive $Y\\in z\\mathcal D(x,q)$, such that (\\ref{6-16-2}) holds.\n\n\n\\end{proof}\n\\begin{Lemma}\\label{lemFinu}\nUnder the conditions of Theorem~\\ref{mainTheorem}, for every $(x,q)\\in\\mathring{\\mathcal K}$ we have\n\\begin{displaymath}\n-\\infty < u(x,q) < \\infty.\n\\end{displaymath}\nand $u<\\infty$ on $\\mathbb R\\times \\mathbb L^\\infty(dK)$.\n\\end{Lemma}\n\\begin{proof}\nLet us fix an arbitrary $(x,q)\\in\\mathring{\\mathcal K}$. Since $\\mathring{\\mathcal K}$ is an open convex cone, there exists $\\lambda\\in(0,1)$, $(x_1, q_1)\\in\\mathring{\\mathcal K}$, and $x_2>0$, such that\n $$(x,q) =\\lambda (x_1, q_1) + (1- \\lambda)(x_2,0).$$\nNote that $(x_2, 0)\\in\\mathring{\\mathcal K}$ by Lemma \\ref{lemma1}. By \\eqref{finValue}, there $c\\in\\mathcal A(x_2,0)$, such that \n\\begin{equation}\\label{2141}\n\\mathbb E\\left[\\int_0^TU\\left(t, (1 - \\lambda) c_t\\right)d\\kappa_t \\right] > -\\infty.\n\\end{equation}\nAs $\\mathcal A(x_1,q_1)\\neq\\emptyset$ (see Remark \\ref{remNonEmpty}), there exists $\\tilde c\\in\\mathcal A(x_1,q_1)$. As $U(t, \\cdot)$ is nondecreasing, we get\n\\begin{displaymath}\nu(x,q) \\geq \\mathbb E\\left[\\int_0^TU(t, \\lambda \\tilde c_t + (1 - \\lambda) c_t)d\\kappa_t \\right] \n\\geq \\mathbb E\\left[\\int_0^TU(t, (1 - \\lambda) c_t)d\\kappa_t \\right] \n> -\\infty,\n\\end{displaymath}\nwhere the last inequality follows from \\eqref{2141}. This implies finiteness of $u$ on $\\mathring{\\mathcal K}$ from below. \n\nIn order to show finiteness from above, let us fix a process $c\\in\\mathcal A(x,q)$, such~that \n\\begin{equation}\\nonumbe\n\\mathbb E\\left[\\int_0^TU(t, c_t)d\\kappa_t \\right] \n> -\\infty.\n\\end{equation}\nBy Lemma \\ref{finitenessOverZ'}, there exists $Y\\in\\mathcal Z'$, such that \n\\begin{equation}\\label{4103}\n\\mathbb E\\left[\\int_0^TV(t,Y_t)d\\kappa_t \\right]<\\infty.\n\\end{equation\nIt follows from Lemma \\ref{6-16-1} that $Y\\in\\mathcal Y(1,\\rho)$ for some $(1,\\rho)\\in \\mathcal L$. Therefore, by Proposition \\ref{prop1}, we get\n\\begin{equation}\\label{4101}\n\\begin{array}{rcl}\n\\mathbb E\\left[\\int_0^T U (s, c_s)d\\kappa_s \\right]& \\leq& \\mathbb E\\left[\\int_0^T U(s,c_s)d\\kappa_s \\right] + x + \\int_0^T\\rho_sq_sdK_s - \\mathbb E\\left[\\int_0^T c_sY_sd\\kappa_s \\right]\\\\\n&\\leq& \\mathbb E\\left[\\int_0^T V(s,Y_s)d\\kappa_s \\right] + x + \\int_0^T\\rho_sq_sdK_s. \\\\\n\\end{array}\n\\end{equation}\nAs $Y$ satisfies (\\ref{4103}), we conclude that $u(x,q) < \\infty$. \nMoreover, for $(x,q)\\in\\mathcal K$, as $\\mathcal A(x,q)\\neq\\emptyset$ by Remark \\ref{remNonEmpty}, every $c\\in\\mathcal A(x,q)$ satisfies \\eqref{4101} (with the same $Y$). \nThis implies that $u<\\infty$ on $\\mathcal K$ and therefore, by \\eqref{3313}, on $\\mathbb R\\times \\mathbb L^\\infty(dK)$. This completes the proof of the lemma.\n\\end{proof}\nWe recall that, for every $(x,q)\\in\\mathcal K$, $\\mathcal D(x,q)$ is defined in \\eqref{auxiliarySets}. Let $cl \\mathcal D(x,q)$ denote the closure of $\\mathcal D(x,q)$ in $\\mathbb L^0(d\\kappa \\times \\mathbb P).$ The following lemma proves a delicate point that, for $(x,q)\\in\\mathring{\\mathcal K}$, by passing from $\\mathcal D(x,q)$ to $cl\\mathcal D(x,q)$, we do not change the auxiliary dual value function. \n\\begin{Lemma}\\label{auxiliaryLemma}\nLet the conditions of Theorem~\\ref{mainTheorem} hold and $(x,q)\\in\\mathring{\\mathcal K}$. Then, for every $z>0$, we have\n\\begin{equation}\\label{2143}\n-\\infty<\\inf\\limits_{Y \\in\\cl \\mathcal D(x,q)}\\mathbb E\\left[\\int_0^TV(s, zY_s)d\\kappa_s \\right] =\n\\inf\\limits_{Y \\in \\mathcal D(x,q)}\\mathbb E\\left[\\int_0^TV(s, zY_s)d\\kappa_s \\right]<\\infty.\n\\end{equation}\n\\end{Lemma}\n\\begin{proof}\nFiniteness from above follows from Lemma \\ref{6-12-1}. To show finiteness of both infima in \\eqref{2143} from below, by Lemma~\\ref{lemFinu} we deduce\nthe existence of $c\\in\\mathcal A(x,q)$, such that\n\\begin{equation}\\label{6-18-2}\n\\mathbb E\\left[\\int_0^TU(s, c_s)d\\kappa_s \\right] > -\\infty.\n\\end{equation}\nLet $Y\\in{\\cl \\mathcal D(x,q)}$ and let $Y^n\\in\\mathcal Y(y^n,r^n)$, $n\\geq 1$, be a sequence in $\\mathcal D(x,q)$ that converges to $Y$ in $\\mathbb L^0(d\\kappa\\times\\mathbb P)$. By Fatou's lemma, Proposition~\\ref{prop1}, and the definition of the set $\\mathcal B(x,q)$ in \\eqref{auxiliarySets}, we get\n\\begin{displaymath}\n\\mathbb E\\left[\\int_0^TY_sc_sd\\kappa_s \\right] \\leq \n\\liminf\\limits_{n\\to\\infty}\\mathbb E\\left[\\int_0^TY^n_sc_sd\\kappa_s \\right] \\leq\n\\sup\\limits_{n\\geq 1}\\left(xy^n + \\int_0^Tr^n_sq_sdK_s\\right) \\leq 1. \n\\end{displaymath}\nTherefore, we obtain\n\\begin{displaymath}\n\\begin{array}{rcl}\n\\mathbb E\\left[ \\int_0^TU(s, c_s)d\\kappa_s\\right]&\\leq& \\mathbb E\\left[ \\int_0^TU(s, c_s)d\\kappa_s\\right] + 1 - \\mathbb E\\left[\\int_0^TY_sc_sd\\kappa_s \\right] \\\\\n&\\leq& \\mathbb E\\left[ \\int_0^TV(s, Y_s)d\\kappa_s\\right] + 1,\\\\\n\\end{array}\n\\end{displaymath}\nwhich together with (\\ref{6-18-2}) implies finiteness of both infima in \\eqref{2143} from below.\n\n\nLet us show equality of two infima in \\eqref{2143}. It follows from Lemma\n\\ref{6-12-1} that for every $z> 0$ there exists a process $Y\\in z\\mathcal D(x,q)$, such that\n$$\\mathbb E\\left[\\int_0^TV(s, Y_s)d\\kappa_s \\right] <\\infty.$$\nLet us fix $z>0$ and let $\\bar Y\\in\\cl\\mathcal D(x,q)$. Also, let $(Y^n)_{n\\in\\mathbb N}$ be a sequence in $\\mathcal D(x,q)$ that converges to $\\bar Y$ $(d\\kappa\\times\\mathbb P)$-a.e. \nLet us fix $\\delta >0$, then by Lemma \\ref{finitenessOverZ'}, there exists $Z'\\in\\mathcal Z'$, such that\n$$\\mathbb E\\left[\\int_0^TV(t, \\delta\\bar y(x,q) Z'_t)d\\kappa_t \\right]<\\infty,$$\nwhere $\\bar y(x,q)$ is defined in \\eqref{defBary}. Note that $\\bar y(x,q) Z'\\in \\mathcal D(x,q)$ by Lemma \\ref{6-12-1} (see \\eqref{6-18-1}). \nTherefore, using Fatou's lemma and monotonicity of $V$ in the spatial variable, we obtain\n\\begin{displaymath}\n\\begin{array}{rcl}\n\\inf\\limits_{Y \\in \\mathcal D(x,q)}\\mathbb E\\left[\\int_0^TV(t, (z+\\delta)Y_s)d\\kappa_s \\right]&\\leq&\n\\limsup\\limits_{n\\to\\infty}\\mathbb E\\left[\\int_0^TV(t, zY^n_s+\\delta \\bar y(x,q)Z'_s)d\\kappa_s \\right]\\\\ \n&\\leq &\\mathbb E\\left[\\int_0^TV(t, z\\bar Y_s + \\delta\\bar y(x,q) Z'_s)d\\kappa_s \\right] \\\\\n&\\leq &\\mathbb E\\left[\\int_0^TV(t, z\\bar Y_s )d\\kappa_s \\right].\\\\\n\\end{array}\n\\end{displaymath}\nTaking the infimum over $Y\\in\\cl\\mathcal D(x,y)$, we deduce that\n\\begin{equation}\\label{21410}\n \\inf\\limits_{Y \\in \\mathcal D(x,q)}\\mathbb E\\left[\\int_0^TV(t, (z+\\delta)Y_s)d\\kappa_s \\right] \\leq \\inf\\limits_{Y\\in\\cl\\mathcal D(x,y)}\\mathbb E\\left[\\int_0^TV(t, zY_s )d\\kappa_s \\right].\n\\end{equation}\nLet us consider \n$$\\phi(z){:=} \\inf\\limits_{Y \\in \\mathcal D(x,q)}\\mathbb E\\left[\\int_0^TV(t, zY_s)d\\kappa_s \\right], \\quad z>0.$$\nBy the first part of the proof (finiteness of both infima), $\\phi$ is finite-valued on $(0,\\infty)$. Convexity of $V$ in the spatial variable implies that $\\phi$ is also convex. Therefore, $\\phi$ is continuous. As \\eqref{21410} holds for every $\\delta>0$, by taking the limit as $\\delta \\downarrow 0$ in \\eqref{21410}, we conclude that both infima in \\eqref{2143} are equal. This completes the proof of the lemma.\n\n\\end{proof}\n\nLet us define\n\\begin{equation}\\label{defE}\n\\mathcal E {:=} \\{(y,r)\\in \\mathcal L: ~v(y,r) < \\infty \\}.\n\\end{equation}\n\\begin{Lemma}\\label{lemFinv}\nUnder the conditions of Theorem~\\ref{mainTheorem}, for every $(y,r)\\in \\mathcal L$, we have\n$$v(y,r) > -\\infty.$$\nTherefore, $v>-\\infty$ on $\\mathbb R\\times\\mathbb L^1(dK)$. The set \n $\\mathcal E$ is a nonempty convex subset of $\\mathcal L$, whose closure in $\\mathbb R\\times\\mathbb L^1(dK)$ equals to $\\mathcal L$, and such that\n\\begin{equation}\\label{6-17-2}\n\\mathcal E= \\bigcup\\limits_{\\lambda \\geq 1 }\\lambda \\mathcal E.\n\\end{equation}\n\n\\end{Lemma}\n\\begin{proof}\nLet us fix $(y,r)\\in \\mathcal L$, then finiteness of $v(y,r)$ from below follows from \\eqref{finValue} and Lemma~\\ref{6-17-1}.\nTo establish the properties of $\\mathcal E$, we observe that the convexity of $\\mathcal E$ and (\\ref{6-17-2}) follow from convexity and monotonicity of $V$, respectively. \n\n\nIn remains to show that the closure of $\\mathcal E$ in $\\mathbb R\\times\\mathbb L^1(dK)$ contains the origin. \nIn \\eqref{auxiliarySets}, let us consider $(x,q) = (1,0)\\in\\mathcal K$. In this case, we have \n $$\\mathcal D(1,0) = \\bigcup\\limits_{(y,r)\\in\\mathcal L: y\\leq 1}\\mathcal Y(y,r)\\subseteq \\mathcal Y(1),$$\nwhere the last inclusion follows from the very definition of $\\mathcal Y(y,r)$'s in \\eqref{defY}.\nAs, by \\eqref{oldY}, $\\mathcal Y(1)$ is closed in $\\mathbb L^0(d\\kappa\\times\\mathbb P)$ and $\\mathcal D(1,0)\\subseteq \\mathcal Y(1)$, we deduce that \n\\begin{equation}\\label{2146}\n\\cl \\mathcal D(1,0)\\subseteq \\mathcal Y(1).\n\\end{equation}\n By Lemma \\ref{6-12-1}, $\\mathcal Z'\\subset \\mathcal D(1,0)$, as $\\bar y(1,0) = 1$. Therefore, by the bipolar theorem of Brannath and Schachermayer, \\cite[Theorem 1.3]{BranSchach}, we get\n\\begin{equation}\\label{2147}\n(\\mathcal Z')^{oo} \\subseteq\\cl\\mathcal D(1,0).\n\\end{equation}\nOn the other hand, Lemma \\ref{finitenessOverZ'} asserts that \n\\begin{equation}\\label{2148}\n(\\mathcal Z')^{oo} = \\mathcal Y(1).\n\\end{equation}\nCombining \\eqref{2146}, \\eqref{2147}, and \\eqref{2148}, we conclude\n\\begin{equation}\\nonumbe\n\\cl \\mathcal D(1,0)= \\mathcal Y(1).\n\\end{equation}\nTherefore, the sets $\\cl\\mathcal D(1,0)= \\mathcal Y(1)$ and $\\mathcal A(1,0)$ satisfy the precise technical assumptions of \\cite[Theorem 3.2]{MostovyiNec}, which, for every $x>0$, grants the existence of $\\widehat c(x)\\in\\mathcal A(x,0)$, the unique maximizer to $w(x)$, where $w$ is defined in \\eqref{defw}.\nFor every $x>0$, we set\n$\nY_{\\cdot}(x){:=} U'({\\cdot}, \\widehat c_{\\cdot}(x)), \\quad \n(d\\kappa\\times\\mathbb P)-a.e\n$$\nBy \\cite[Theorem 3.2]{MostovyiNec}, for every $x>0$, $Y(x)$ satisfies\n$$Y(x) \\in w'(x)\\cl\\mathcal D(1,0)\\quad and\\quad \\mathbb E\\left[ \\int_0^TV(t, Y_t(x))d\\kappa_t\\right]<\\infty,\\quad x>0,$$\nwhere by \\cite[Theorem 3.2]{MostovyiNec}, $w$ is a strictly concave, differentiable function on $(0,\\infty)$ that satisfies the Inada conditions. \nTherefore, as $w'(x)$ can be arbitrary close to $0$ (by taking $x$ large enough an by using the Inada conditions) and by Lemmas \\ref{6-12-1} and \\ref{auxiliaryLemma}, we conclude that the closure of $\\mathcal E$ in $\\mathbb R\\times \\mathbb L^1(dK)$ contains origin. \n\nIn order to prove that the closure of $\\mathcal E$ in $\\mathbb R\\times \\mathbb L^1(dK)$ equals to $\\mathcal L$, let $(y,r)\\in\\mathcal L\\backslash (0, 0)$ be fixed. Let us take $\\varepsilon>0$. We want to find $(\\tilde y, \\tilde r)$, such that \n\\begin{equation}\n\\label{4171}\n|\\tilde y - y| + || \\tilde r- r ||_{\\mathbb L^1(dK)}<\\varepsilon,\n\\end{equation}\nand\n\\begin{equation}\\label{4172}\n(\\tilde y, \\tilde r)\\in\\mathcal E.\n\\end{equation}\nAs the closure of $\\mathcal E$ in $\\mathbb R\\times \\mathbb L^1(dK)$ contains origin, we can pick $(y^0, r^0)\\in\\mathcal E$, such that \n\\begin{equation}\\label{4173}\n|y^0| + ||r^0||_{\\mathbb L^1(dK)} \\leq \\varepsilon\/3\n\\end{equation}\nand $Y\\in\\mathcal Y(y^0,r^0)$, such that \n\\begin{equation}\\label{4174}\n\\mathbb E\\left[\\int_0^T V(t, Y_t)d\\kappa_t \\right]<\\infty.\n\\end{equation}\nLet us fix $\\alpha >1$, such that \n\\begin{equation}\\label{4177}\n\\frac{|y| + ||r||_{\\mathbb L^1(dK)}}{\\alpha} \\leq \\varepsilon\/3\n\\end{equation}\n and set $\\varepsilon' {:=} \\frac{1}{\\alpha}\\in(0,1)$. By \\eqref{6-17-2}, $(\\alpha y^0, \\alpha r^0)\\in\\mathcal E.$ Let \n\\begin{equation}\\nonumbe\n\\tilde y {:=} (1-\\varepsilon')y + \\varepsilon'\\alpha y^0,\\quad \\tilde r {:=} (1-\\varepsilon')r + \\varepsilon'\\alpha r^0.\n\\end{equation} \nThen\n\\begin{displaymath}\n\\begin{array}{rcl}\n|y - \\tilde y| + || r - \\tilde r||_{\\mathbb L^1(dK)} &= & \\varepsilon' \\alpha |\\frac{y}{\\alpha} - y^0| + \\varepsilon'\\alpha || \\frac{r}{\\alpha} - r^0||_{\\mathbb L^1(dK)}\\\\\n&\\leq & \\frac{|y| + ||r||_{\\mathbb L^1(dK)}}{\\alpha} + |y^0| + ||r^0||_{\\mathbb L^1(dK)}\n\\\\\n&\\leq & \\frac{2\\varepsilon}{3},\n\\end{array}\n\\end{displaymath}\nwhere in the last inequality we have used \\eqref{4173} and \\eqref{4177}. Thus $(\\tilde y, \\tilde r)$ satisfies \\eqref{4171}. Further, as $0\\in\\mathcal Y(y,r)$, by convexity of $\\mathcal L$ and using Proposition \\ref{prop1}, we get\n\\begin{displaymath}\nY=(1-\\varepsilon') 0 + \\varepsilon' \\alpha Y \\in\\mathcal Y(\\tilde y, \\tilde r),\n\\end{displaymath}\nwhich by \\eqref{4174} implies \\eqref{4172}. This completes the proof of the lemma.\n\\end{proof}\n\n\\subsection{Existence and uniqueness of solutions to \\eqref{primalProblem} and \\eqref{dualProblem}; semicontinuity and biconjugacy of $u$ and $v$}\n\n\\begin{Lemma}\\label{existenceUniqueness}\nUnder the conditions of Theorem~\\ref{mainTheorem},\nthe value function $v$ is convex, proper, and lower semicontinuous with respect to the topology of $\\mathbb R\\times \\mathbb L^1(dK)$. For every $(y,r)\\in\\mathcal E$, there exists a unique solution to \n(\\ref{dualProblem}). Likewise, $u$ is concave, proper, and upper semicontinuous with respect to the strong topology of $\\mathbb R\\times \\mathbb L^\\infty(dK)$. For every $(x,q)\\in\\{u>-\\infty\\}$ there exists a unique solution to~(\\ref{primalProblem}).\n\\end{Lemma}\n\\begin{proof}\nLet $(y^n, r^n)_{n\\in\\mathbb N}$ be a sequence in $\\mathcal L$ that converges to $(y,r)$ in $\\mathbb R\\times\\mathbb L^1(dK)$. \nPassing if necessary to a subsequence, we will assume that \n\\begin{equation}\\label{2151}\n\\lim\\limits_{n\\to\\infty} v(y^n,r^n)= \\liminf\\limits_{n\\to\\infty}v(y^n, r^n).\n\\end{equation}\nLet $Y^n\\in\\mathcal Y(y^n, r^n)$, $n\\in\\mathbb N$, be such that \n\\begin{equation}\\label{2152}\n\\mathbb E\\left[ \\int_0^TV(t,Y^n_t)d\\kappa_t\\right]\\leq v(y^n, r^n) +\\frac{1}{n}, \\quad n\\in\\mathbb N.\n\\end{equation}\nBy passing to convex combinations and applying Komlos'-type lemma, see e.g. \\cite[Lemma A1.1]{DS}, we may suppose that $\\widetilde Y^n\\in\\conv\\left(Y^n, Y^{n+1},\\dots\\right)$, $n\\in\\mathbb N$, converges $(d\\kappa\\times\\mathbb P)$-a.e. to some $\\widehat Y$. \n\nFor every $(x,q)\\in\\mathcal K$ and $c\\in\\mathcal A(x,q)$, by Fatou's lemma, we have\n\\begin{displaymath}\n\\begin{array}{c}\n\\mathbb E\\left[ \\int_0^T c_t\\widehat Y_td\\kappa_t\\right]\\leq \\liminf\\limits_{n\\to\\infty}\\mathbb E\\left[ \\int_0^T c_t \\widetilde Y^n_td\\kappa_t\\right]\\leq\nxy + \\int_0^Tq_sr_sdK_s.\n\\end{array}\n\\end{displaymath} \nTherefore, by Proposition \\ref{prop1}, $\\widehat Y\\in\\mathcal Y(y,r)$. With $\\bar y {:=} \\sup\\limits_{n\\geq 1}y^n$, we have $(\\widetilde Y^n)_{n\\in\\mathbb N}\\subseteq \\mathcal Y(\\bar y).$ Therefore, by \\cite[Lemma 3.5]{MostovyiNec}, we deduce that \n$V^{-}(t,\\widetilde Y^n_t)$, $n\\in\\mathbb N$, is a uniformly integrable sequence. \nCombining uniform integrability with the convexity of $V$ in the spatial variable, we get\n\\begin{equation}\\label{4141}\n\\begin{array}{c}\nv(y,r)\\leq \\mathbb E\\left[ \\int_0^TV(t,\\widehat Y_t)d\\kappa_t\\right]\\leq \n\\liminf\\limits_{n\\to\\infty}\n\\mathbb E\\left[ \\int_0^TV(t, \\widetilde Y^n_t)d\\kappa_t\\right]\n\\\\\n\\leq \n\\liminf\\limits_{n\\to\\infty}\n\\mathbb E\\left[ \\int_0^TV(t, Y^n_t)d\\kappa_t\\right] = \\liminf\\limits_{n\\to\\infty}v(y^n,r^n),\n\\end{array}\n\\end{equation}\nwhere in the last equality we have used \\eqref{2151} and \\eqref{2152}. Since $(y^n, q^n)$ was an arbitrary sequence that converges to $(y,r)$, lower semicontinuity of $v$ in strong topology of $\\mathbb R\\times \\mathbb L^1(dK)$ follows. Since $\\mathcal L$ is closed and $v = \\infty$ outside of $\\mathcal L$, we deduce that $v$ is lower semicontinuous on $\\mathbb R\\times \\mathbb L^1(dK)$. The function $v$ is proper by Lemma \\ref{lemFinv}. Note that \\eqref{4141} also implies that $\\mathcal E$ defined in \\eqref{defE} is $\\mathbb R\\times\\mathbb L^1(dK)$-norm closed. For $(y,r)\\in\\mathcal E$, by taking $(y^n, r^n)= (y,r)$, $n\\in\\mathbb N$, we deduce the existence of a minimizer to \\eqref{dualProblem}. Strict convexity of $V$ results in the uniqueness of the minimizer to \\eqref{dualProblem}. Convexity of $v$ follows. Upper semicontinuity of $u$ with respect to the norm-topology of $\\mathbb R\\times \\mathbb L^{\\infty}(dK)$ can be proven similarly, first proving semi-continuity on \n $\\mathcal K$ by a Fatou-type argument, then using the closedness of $\\mathcal{K}$ and the definition of $u$ outside it. \n\\end{proof}\n\\begin{Corollary}\nUnder the conditions of Theorem~\\ref{mainTheorem},\n$-u$ and $v$ are also lower semicontinuous with respect to the weak topologies $\\sigma(\\mathbb R\\times\\mathbb L^{\\infty}(dK), (\\mathbb R\\times\\mathbb L^{\\infty})^*(dK))$ and $\\sigma(\\mathbb R\\times\\mathbb L^1(dK),\\mathbb R\\times\\mathbb L^{\\infty}(dK))$, respectively.\n\\end{Corollary}\n\\begin{proof}\nThe assertions of the corollary is a consequence of \\cite[Proposition 2.2.10]{BarbuPrec}, see also \\cite[Corollary I.2.2]{EkelandTemam}.\n\\end{proof}\n\nWe recall that ${\\cl \\mathcal D(x,q)}$ denotes the closure of $\\mathcal D(x,q)$ in $\\mathbb L^0(d\\kappa\\times\\mathbb P)$. \n\\begin{Lemma}\\label{5181}\nUnder the conditions of Theorem \\ref{mainTheorem}, for every $(x,q)$ in $\\mathring{\\mathcal K}$, and a nonnegative optional process $c$, we have\n$$\nc\\in\\mathcal A(x,q)\\quad{if~and~only~if}\\quad \\mathbb E\\left[\\int_0^T c_sY_sd\\kappa_s \\right]\\leq 1\\quad {for~every~}Y\\in{\\cl \\mathcal D(x,q)}.\n$$\n\\end{Lemma}\n\\begin{proof}\nLet $(x,q)$ in $\\mathring{\\mathcal K}$, $c$ is a nonnegative optional process such that \n\\begin{equation}\\label{5191}\n\\mathbb E\\left[ \\int_0^T c_sY_sd\\kappa_s \\right]\\leq 1\\quad {for~every}\\quad Y\\in{\\cl \\mathcal D(x,q)}.\n\\end{equation}\nConsider arbitrary $Z\\in\\mathcal Z'$ and $\\Lambda\\in\\Upsilon$. Let the corresponding $r$ be given by \\eqref{6-25-1} and we set $$y'{:=} x + \\int_0^Tq_sr_sdK_s.$$\n\nIf $y' = 0$, then $y\\Lambda Z\\in\\mathcal D(x,q)$ for every $y>0$. Thus, by \\eqref{5191}, we obtain that $\\mathbb E\\left[\\int_0^T c_sy\\Lambda_sZ_sd\\kappa_s \\right]\\leq 1$. Taking the limit as $y\\to\\infty$, we get\n\\begin{equation}\\label{5183}\\mathbb E\\left[\\int_0^T c_s\\Lambda_sZ_sd\\kappa_s \\right] = 0= x + \\int_0^Tq_sr_sdK_s= x + \\mathbb E\\left[\\int_0^Tq_se_s\\Lambda_sZ_sd\\kappa_s \\right],\n\\end{equation}\nwhere in the last equality we have used \\eqref{6-25-1}. \n\nIf $y'>0$, \nthen $\\frac{1}{y'}\\Lambda Z\\in\\mathcal D(x,q)$ and thus by \\eqref{5191}, we obtain $$\\mathbb E\\left[ \\int_0^T c_s\\frac{1}{y'}\\Lambda_sZ_sd\\kappa_s\\right]\\leq 1 = \\frac{ x + \\int_0^Tq_sr_sdK_s}{y'} = \\frac{1}{y'}\\left( x + \\mathbb E\\left[\\int_0^T q_se_s\\Lambda_sZ_sd\\kappa_s\\right]\\right),$$\nwhere in the last equality, we have used \\eqref{6-25-1} again. Consequently, we deduce\n$$\\mathbb E\\left[ \\int_0^T c_s\\Lambda_sZ_sd\\kappa_s\\right]\\leq x + \\mathbb E\\left[\\int_0^T q_se_s\\Lambda_sZ_sd\\kappa_s\\right],$$\nwhich together with \\eqref{5183}, by Lemma \\ref{9-5-1}, imply that $c\\in\\mathcal A(x,q)$.\n\nConversely, let $(x,q)\\in\\mathring{\\mathcal K}$, $c\\in\\mathcal A(x,q)$ and $Y\\in {\\rm cl}\\mathcal D(x,q)$. Then there exists a sequence $Y^n\\in\\mathcal Y(y^n, r^n)$ convergent to $Y$, $(d\\kappa\\times\\mathbb P)$-a.e., where $(y^n,r^n)\\in\\mathcal B(x,q)$. As, \n$$\\mathbb E\\left[\\int_0^T c_sY^n_sd\\kappa_s \\right]\\leq 1,\\quad n\\in\\mathbb N,$$\nby Fatou's lemma, we get\n$$\\mathbb E\\left[\\int_0^T c_sY_sd\\kappa_s \\right]\\leq \\liminf\\limits_{n\\to\\infty}\\mathbb E\\left[\\int_0^T c_sY^n_sd\\kappa_s \\right]\\leq 1.$$\nThis completes the proof of the lemma.\n\\end{proof}\n\\begin{Lemma}\\label{conjugacy}\nUnder the conditions of Theorem \\ref{mainTheorem}, for every $(x,q)$ in $\\mathring{\\mathcal K}$, we have\n\\begin{equation}\\label{conjugacyOneDirection}\nu(x,q) = \\inf\\limits_{(y,r)\\in \\mathcal L}\\left(v(y,r) + xy + \\int_0^Tr_sq_sdK_s\\right).\n\\end{equation}\n\\end{Lemma}\n\\begin{proof} Let us fix $(x,q)\\in\\mathring{\\mathcal K}$. \nBy Lemma \\ref{6-12-1}, $\\mathcal A(x,q)$ and ${\\cl \\mathcal D(x,q)}$ contain strictly positive elements. \nTherefore, using Lemma \\ref{5181}\nwe deduce that the sets $\\mathcal A(x,q)$ and ${\\cl \\mathcal D(x,q)}$ satisfy the assumptions \nof \\cite[Theorem 3.2]{MostovyiNec}. From this theorem, Lemma~\\ref{auxiliaryLemma}, and the definition of the set $\\mathcal B(x,q)$, we get\n\\begin{equation}\\nonumbe\n\\begin{array}{rcl}\nu(x,q) &=& \\inf\\limits_{z>0}\\left(\\inf\\limits_{Y\\in{\\cl \\mathcal D(x,q)}}\\mathbb E\\left[ \\int_0^TV(t, zY_s)d\\kappa_s\\right] + z \\right) \\\\\n &=& \\inf\\limits_{z>0}\\left(\\inf\\limits_{Y\\in\\mathcal D(x,q)}\\mathbb E\\left[ \\int_0^TV(t, zY_s)d\\kappa_s\\right] + z \\right) \\\\\n&=& \\inf\\limits_{z>0}\\left(\\inf\\limits_{(y,r)\\in z\\mathcal B(x,q)}v(y,r) + z \\right) \\\\\n&\\geq& \\inf\\limits_{(y,r)\\in \\mathcal L}\\left(v(y,r) + xy + \\int_0^Tq_sr_sdK_s \\right).\\\\\n\\end{array}\n\\end{equation}\nCombining this with the conclusion of Lemma~\\ref{6-17-1}, we deduce that \\eqref{conjugacyOneDirection} holds for every $(x,q)\\in\\mathring{\\mathcal K}$. \n\n\n\\end{proof}\n\nBefore proving the biconjugacy relations of item $(iii)$, Theorem \\ref{mainTheorem}, we need a preliminary lemma. \nEssentially following the notations in \\cite{EkelandTemam}, we define \n\\begin{equation}\\label{v*}\nv^*(x,q){:=} \\inf\\limits_{(y,r)\\in\\mathbb R\\times \\mathbb L^1(dK)}\\left(v(y,r) + xy + \\int_0^Tq_sr_sdK_s \\right),\\quad (x,q)\\in\\mathbb R\\times \\mathbb L^{\\infty}(dK).\n\\end{equation}\n\\begin{equation}\\label{v**}\nv^{**}(y,r){:=} \\sup\\limits_{(x,q)\\in\\mathbb R\\times \\mathbb L^\\infty(dK)}\\left(v^*(x,q) - xy - \\int_0^Tq_sr_sdK_s \\right),\\quad (y,r)\\in\\mathbb R\\times \\mathbb L^1(dK).\n\\end{equation}\n\\begin{Remark}In \\cite{EkelandTemam}, conjugate convex functions are considered on general spaces $V$ and $V^*$ supplied with $\\sigma(V,V^*)$ and $\\sigma(V^*, V)$ topologies, which in our case are $V= \\mathbb R\\times\\mathbb L^1(dK)$, $V^* = \\mathbb R\\times\\mathbb L^\\infty(dK)$. Thus, the starting point of our analysis is $v$, not $u$. We remind the reader we have already proved that the dual value function $v$ is convex, proper and lower-semicontinuous on the space $V= \\mathbb R\\times\\mathbb L^1(dK)$. \n\\end{Remark}\n\\begin{Lemma}\\label{lemv**}\nUnder the conditions of Theorem \\ref{mainTheorem}, we have \n\\begin{equation}\n\\label{2166}\nv^{**} = v, \n\\end{equation}\n\\begin{equation}\n\\label{4151}\nv^*(x,q) = -\\infty, \\quad for~every\\quad (x,q)\\in\\mathbb R\\times\\mathbb L^\\infty(dK)\\backslash \\mathcal K.\n\\end{equation}\n\n\\end{Lemma}\n\n\\begin{proof}\nTo show \\eqref{2166}, we observe that by Lemma \\ref{existenceUniqueness}, $v$ is lower semicontinuous in the $\\mathbb R\\times\\mathbb L^1(dK)$-norm topology (and therefore, by \\cite[Corollary I.2.2]{EkelandTemam}, also in the weak topology $\\sigma(\\mathbb R\\times\\mathbb L^1(dK),\\mathbb R\\times\\mathbb L^\\infty(dK))$). As a result, by \\cite[Proposition I.4.1]{EkelandTemam}, we get~\\eqref{2166}.\n\nThe proof of \\eqref{4151} will be done in several steps. \n\n{\\it Step 1.} Let $(x,q)\\in\\mathbb R\\times\\mathbb L^\\infty(dK)\\backslash \\mathcal K$. According to Proposition \\ref{prop1}, there exists $(y,r)\\in\\mathcal L$, such that\n\\begin{equation}\\nonumbe\nC{:=} xy + \\int_0^T q_sr_sdK_s <0.\n\\end{equation}\nTherefore, as $\\mathcal L$ is a cone, for every $a>0$, $(ay, ar)\\in\\mathcal L$, and we have\n\\begin{equation}\\label{4153}\nxay + \\int_0^Tq_sar_sdK_s = aC<0.\n\\end{equation}\nNote that $0\\in\\mathcal Y(ay, ar)$, $a>0$.\n\n{\\it Step 2.} Let us consider $Z\\in\\mathcal Z'$, such that \n\\begin{equation}\\label{4154}\n\\mathbb E\\left[\\int_0^TV\\left(t,\\tfrac{1}{2}Z_t\\right)d\\kappa_t\\right]<\\infty.\n\\end{equation}\nThe existence of such a $Z$ is granted by Lemma \\ref{finitenessOverZ'}.\nFurther, by Corollary \\ref{cor1}, there exists $\\rho\\in\\mathbb L^1(dK)$, such that $(1, \\rho)\\in\\mathcal L$, $Z\\in\\mathcal Y(1,\\rho)$, and\n\\begin{equation}\\nonumbe\n\\int_0^t\\rho_sdK_s = \\mathbb E\\left[\\int_0^t Z_se_sd\\kappa_s \\right],\\quad t\\in[0,T].\n\\end{equation}\nLet us set \n\\begin{equation}\\nonumbe\nD{:=} x + \\int_0^Tq_s\\rho_sdK_s\\in\\mathbb R.\n\\end{equation}\n\n{\\it Step 3.} In \\eqref{4153}, let us pick \n\\begin{equation}\\nonumbe\na = \\frac{|D| + 1}{-C}>0.\n\\end{equation}\nThen, we have \n\\begin{equation}\\label{4158}\naC + D = -|D| + D -1<0.\n\\end{equation}\n\n{\\it Step 4.} Let us define \n\\begin{equation}\\label{4159}\nY{:=} \\tfrac{1}{2} Z, \\quad y'{:=} \\tfrac{1}{2}ay + \\tfrac{1}{2}, \\quad and \\quad r'{:=} \\tfrac{1}{2}ar + \\tfrac{1}{2}\\rho.\n\\end{equation}\nThen by \\eqref{4154}, we obtain\n\\begin{equation}\\label{41510}\n\\mathbb E\\left[\\int_0^TV\\left(t,Y_t\\right)d\\kappa_t\\right]<\\infty.\n\\end{equation}\nAs $Z\\in\\mathcal Y(1,\\rho)$ and $0\\in\\mathcal Y(ay, ar)$, by convexity of $\\mathcal L$ and Proposition \\ref{prop1}, we have\n\\begin{equation}\\label{41511}\nY\\in\\mathcal Y(y',r'),\n\\end{equation}\nwhere $y'$ and $r'$ are defined in \\eqref{4159}. Now, it follows from \\eqref{41510} and \\eqref{41511} that $(y', r')\\in\\mathcal E$ (where $\\mathcal E$ is defined in \\eqref{defE}). \nTherefore, we obtain\n\\begin{displaymath}\n\\begin{array}{c}\n2(xy' + \\int_0^Tq_sr'_sdK_s) = (ay + 1)x + \\int_0^T(ar_s + \\rho_s)q_sdK_s \\\\= a(xy + \\int_0^Tq_sr_sdK_s) + x + \\int_0^Tq_s\\rho_sdK_s = aC + D <0,\\\\\n\\end{array}\n\\end{displaymath}\nwhere the last inequality follows from \\eqref{4158}.\nTo recapitulate, we have shown the existence of $(y', r')$, such that\n\\begin{equation}\\label{41512}\n(y', r')\\in\\mathcal E\\quad and\\quad xy' + \\int_0^Tq_sr'_sdK_s<0.\n\\end{equation}\n\n{\\it Step 6.} For $y'$ and $r'$ defined in \\eqref{4159}, as $v(y', r') <\\infty$ and $xy' + \\int_0^Tq_sr'_sdK_s<0$ by \\eqref{41512}, from the monotonicity of $V$, we get\n$$\\infty>v(y',r')\\geq v(\\lambda y', \\lambda r'),\\quad \\lambda \\geq 1.$$\nAs $\\bigcup\\limits_{\\lambda\\geq 1}(\\lambda y', \\lambda r')\\subset \\mathcal L$, we conclude via \\eqref{41512} that\n$$v^*(x,q) \\leq \\lim\\limits_{\\lambda \\to\\infty}\\left( v(\\lambda y', \\lambda r') + \\lambda \\left(xy' + \\int_0^Tq_sr'_sdK_s\\right)\\right) = -\\infty.$$\nTherefore, \\eqref{4151} holds. This completes the proof of the lemma.\n\\end{proof}\n\\begin{Lemma}\\label{lemBiconjugacy}\nUnder the conditions of Theorem \\ref{mainTheorem}, we have \n\\begin{equation}\\label{2161\nv(y,r) = \\sup\\limits_{(x,q)\\in\\mathcal K}\\left(u(x,q) - xy - \\int_0^Tr_sq_sdK_s\\right),\\quad (y,r)\\in\\mathcal L,\n\\end{equation}\n\\begin{equation}\\label{21610}\nu(x,q) = \\inf\\limits_{(y,r)\\in\\mathcal L}\\left(v(y,r) + xy + \\int_0^Tq_sr_sdK_s\\right),\\quad (x,q)\\in \\mathcal K.\n\\end{equation}\n\\end{Lemma}\n\\begin{proof}\n\n\nLemma \\ref{conjugacy} and \\eqref{3313} imply that on $\\mathring{\\mathcal K}$, for $v^*$ defined in \\eqref{v*}, we have\n\\begin{equation}\\label{2164}\nv^{*} = u\n\\end{equation}\nBy Lemma \\ref{lemv**}, $v^* = -\\infty$ on $\\mathbb R\\times\\mathbb L^\\infty(dK)\\backslash\\mathcal K$. \nFrom \\cite[Definition I.4.1]{EkelandTemam} and Lemma \\ref{existenceUniqueness}, respectively, we deduce that both $v^*$ and $u$ are upper semicontinuous in the topology of $\\mathbb R\\times\\mathbb L^\\infty(dK)\n. Consequently, from \\eqref{2164}, using \\cite[Corollary I.2.1]{EkelandTemam}, \nwe get\n\\begin{equation}\\label{2201}v^* = u,\\quad on\\quad\\mathbb R\\times\\mathbb L^\\infty(dK).\n\\end{equation}\nAs a result, with $v^{**}$ being defined in \\eqref{v**}, for every $(y,r)\\in\\mathbb R\\times \\mathbb L^1(dK)$, we obtain\n\\begin{equation}\\label{2165}\nu^{*}(y,r){:=} \\sup\\limits_{(x,q)\\in\\mathbb R\\times \\mathbb L^\\infty(dK)}\\left(u(x,q) - xy - \\int_0^Tq_sr_sdK_s \\right)=v^{**}(y,r)\n\\end{equation}\nTherefore, from \\eqref{2166} in Lemma \\ref{lemv**} and \\eqref{2165}, we get\n\\begin{equation}\\nonumbe\nv = u^{*},\\quad on\\quad\\mathbb R\\times\\mathbb L^1(dK).\n\\end{equation}\nAs a result, applying Lemma \\ref{lemv**} again and since $u = -\\infty$ outside of $\\mathcal K$ by \\eqref{3313}, we deduce \n\\begin{equation}\\nonumbe\n\\begin{array}{rcl}\nv(y,r) &=& \\sup\\limits_{(x,q)\\in\\mathbb R\\times\\mathbb L^\\infty(dK)}\\left( u(x,q)- xy - \\int_0^Tr_sq_sdK_s\\right)\\\\\n&=& \\sup\\limits_{(x,q)\\in\\mathcal K}\\left( u(x,q) - xy - \\int_0^Tr_sq_sdK_s\\right),\\quad(y,r)\\in\\mathcal L,\\\\\n\\end{array}\n\\end{equation}\nThus, \\eqref{2161} holds. \n\nIn turn, from \\eqref{2201} using \\eqref{3313}, \nwe conclude that\n\n\\begin{equation}\\nonumber\n\\begin{array}{rcl}\nu(x,q) &=& \\inf\\limits_{(y,r)\\in\\mathbb R\\times \\mathbb L^1(dK)}\\left(v(y,r) + xy + \\int_0^Tq_sr_sdK_s\\right),\\\\\n&=& \\inf\\limits_{(y,r)\\in\\mathcal L}\\left(v(y,r) + xy + \\int_0^Tq_sr_sdK_s\\right),\\quad (x,q)\\in\\mathcal K,\\\\\n\\end{array}\n\\end{equation}\nwhich proves \\eqref{21610} and extends the assertion of Lemma \\ref{conjugacy} to the boundary of $\\mathring{\\mathcal K}$.\n\n\\end{proof}\n\n\\begin{proof}[Proof of Theorem \\ref{mainTheorem}]\nThe assertions of item $(i)$ follow from Lemmas \\ref{lemFinu} and \\ref{lemFinv}, item $(ii)$ results from Lemma \\ref{existenceUniqueness}, whereas the validity of item $(iii)$ come from Lemma\n\\ref{lemBiconjugacy}. This completes the proof of the theorem.\n\\end{proof}\n\n\n\n\\section{Subdifferentiability of $u$}\\label{secSubdifferentiability}\nIn order to establish subdifferentiability of $u$, we need to strengthen \nAssumption~\\ref{asEnd} and to impose the following condition. \n\\begin{Assumption}\\label{asACclock}\nThere exists an a.s. bounded away from $0$ and $\\infty$ \nprocess $\\varphi$, such that\n\\begin{equation}\\nonumbe\nd\\kappa(\\omega) =\\varphi d K, \\quad for\\quad\\mathbb P-a.e.\\quad\\omega \\in\\Omega.\n\\end{equation}\n\n\\end{Assumption}\n\n\n\nLet \n\\begin{equation}\\label{defPi}\n\\begin{array}{rcl}\n\\mathcal P &{:=} &\\left\\{ \\rho:~(1,\\rho)\\in\\mathcal L\\right\\},\\\\\n\\end{array}\n\\end{equation}\n\\begin{Remark} \n\n $\\mathcal P$ defined in \\eqref{defPi} needs to be uniformly integrable with respect to the measure $dK$ in order for the proof of subdifferentiability of $u$ to go through.\n Assumption \\ref{asEnd} through Lemma \\ref{lemma1} only implies \n that $\\mathcal P$ is $\\mathbb{L}^1$ bounded. A stronger condition on the stochastic clock and income stream is, therefore, needed to obtain uniform integrability. \n \n\n\\end{Remark}\nThe following theorem characterizes subdifferentiability of $u$ over $\\mathring{\\mathcal K}$,\nwhere we are looking for an $\\mathbb R\\times\\mathbb L^1(dK)$-valued subgradient. Under our assumptions, we can find elements of the sub gradient which both belong to the effective domain of $v$, $\\mathcal E$, and are bounded, i.e. in $\\mathbb{R}\\times \\mathbb{L}^{\\infty} (dK)$.\n\\begin{Theorem}\\label{mainTheorem2} \nLet the conditions of Theorem~\\ref{mainTheorem} and Assumption \\ref{asACclock} hold. Then for every $(x,q)\\in\\mathring{\\mathcal K}$, the subdifferential of $u$ at $(x,q)$ is a nonempty and contains an element of $\\mathcal E$, i.e.,\n\\begin{equation}\\label{subdifNonempty}\n\\partial u(x,q)\\cap \\mathcal E\\neq \\emptyset.\n\\end{equation}\nMoreover, for $(x,q)\\in\\mathring{\\mathcal K}$ and $(y,r)\\in\\mathcal L$, $(y,r)\\in\\partial u(x,q)$ if and only if the following conditions hold:\n\\begin{equation}\\label{2191}\n |v(y,r)|<\\infty,\n \\end{equation}\n thus, $(y,r)\\in\\mathcal E$, \n \\begin{equation}\\label{2192}\n \\mathbb E\\left[\\int_0^T\\widehat Y_t(y,r) \\widehat c_t(x,q)d\\kappa_t\\right] = xy + \\int_0^Tq_sr_sdK_s,\n \\end{equation}\n \\begin{equation}\\label{2193}\n \\widehat Y_t(y,r) = U'(t,\\widehat c_t(x,q)), \\quad (d\\kappa\\times\\mathbb P)- a.e.,\n \\end{equation}\n where $\\widehat c(x,q)$ and $\\widehat Y(y,r)$ are the unique optimizers to \\eqref{primalProblem} and \\eqref{dualProblem}, respectively.\n\\end{Theorem}\n\n\n\n\n\n\\subsection{Uniform integrability of $\\mathcal P$}\n\\begin{Lemma}\\label{uiP}\nLet the conditions of Theorem~\\ref{mainTheorem2} hold. Then $\\mathcal P$ is $\\mathbb{L}^{\\infty}(dK)$-bounded, and, therefore, a uniformly integrable family.\n\\end{Lemma}\n\\begin{proof}\n{\\it Step 1.} For an arbitrary $q:~[0,T] \\to [0,1]$, let us define\n\\begin{displaymath}\n\\beta(q) {:=} \\sup\\limits_{\\mathbb Q\\in\\mathcal M} \\mathbb {E^Q}\\left[ \\int_0^T q(s) |e_s|d\\kappa_s\\right] = \\sup\\limits_{\\mathbb Q\\in\\mathcal M, \\tau\\in\\Theta} \\mathbb {E^Q}\\left[ \\int_0^{\\tau} q(s) |e_s|d\\kappa_s\\right].\n\\end{displaymath}\nNote that by Assumption \\eqref{asEnd}, \nwe have\n\\begin{equation}\\nonumber\n\\begin{array}{rcl}\n\\beta(q) &=&\\sup\\limits_{\\mathbb Q\\in\\mathcal M} \\mathbb {E^Q}\\left[ \\int_0^T q(s) |e_s|d\\kappa_s\\right]\\\\\n&\\leq&\\sup\\limits_{\\mathbb Q\\in\\mathcal M} \\mathbb {E^Q}\\left[ \\int_0^T q(s) X'_s\\varphi d K_s\\right]\\\\\n&\\leq&\\int_0^T q(s)\\sup\\limits_{\\mathbb Q\\in\\mathcal M} \\mathbb {E^Q}\\left[ CX'_s\\right]d K_s\\\\\n&\\leq&CX'_0\\int_0^T q(s)d K_s,\\\\\n\\end{array}\n\\end{equation}\nwhere $C\\in\\mathbb R$ is such that $|\\varphi| \\leq C$. \nIt follows from \\cite[Proposition 4.3]{K96} and \\cite[Theorem 3.1]{FK} that there exists a nonnegative c\\`adl\\`ag process $V = \\beta(q) + H\\cdot S$, such that\n\\begin{displaymath}\nV_t \\geq \\esssup\\limits_{\\mathbb Q\\in \\mathcal M, \\tau \\in\\Theta, \\tau\\geq t}\\mathbb {E^Q}\\left[\\int_0^{\\tau}q(s)|e_s|d\\kappa_s | \\mathcal F_t \\right],\\quad t\\in[0,T].\n\\end{displaymath} \nConsequently, $V$ satisfies\n$$V_\\tau \\geq \\int_0^\\tau q(s)|e_s|d\\kappa_s \\geq \\int_0^\\tau q(s)e_sd\\kappa_s,\\quad \\Pas,\\quad\\tau \\in\\Theta,$$\nand thus, $(\\beta(q), -q)\\in\\mathcal K$. As a result, from the definitions of $\\mathcal L$ and $\\mathcal P$, we get\n\\begin{equation}\\nonumbe\n\\beta(q)\\geq \\sup\\limits_{\\rho \\in\\mathcal P}\\int_0^T q(s)\\rho(s)dK_s.\n\\end{equation}\nOne can see that\n for every $\\rho\\in\\mathcal P$, we have\n$\\rho \\leq f{:=} {CX'_0}$, $d K$-a.e. \n\n\n{\\it Step 2.} For an arbitrary $q:~[0,T] \\to [-1,0]$, let us set \n\\begin{displaymath}\n\\tilde \\beta(q) {:=} \\sup\\limits_{\\mathbb Q\\in\\mathcal M}\\mathbb {E^Q}\\left[ \\int_0^T q(s) (-|e_s|) d\\kappa_s \\right].\n\\end{displaymath}\nAs in {\\it Step 1}, we can construct a c\\`adl\\`ag process $\\tilde V = \\tilde \\beta(q) + H\\cdot S$, s.t.\n\\begin{displaymath}\n\\tilde V_t \\geq \\esssup\\limits_{\\mathbb Q\\in\\mathcal M, \\tau \\in\\Theta, \\tau \\geq t}\\mathbb {E^Q}\\left[ \\int_0^T q(s)(-|e_s|) d\\kappa_s |\\mathcal F_t\\right] \\geq \\esssup\\limits_{\\mathbb Q\\in\\mathcal M, \\tau \\in\\Theta, \\tau \\geq t}\\mathbb {E^Q}\\left[ \\int_0^T q(s)e_s d\\kappa_s |\\mathcal F_t\\right].\n\\end{displaymath}\nThis implies that $(\\tilde \\beta(q), -q)\\in\\cl\\mathcal K$. Therefore, \n\\begin{displaymath}\n\\beta(q) \\geq \\int_0^T q(s)\\rho (s) dK_s,\\quad for~every~\\rho \\in\\mathcal P.\n\\end{displaymath}\nSimilarly to {\\it Step 1},\none can see that $\\rho \\geq -f$, $d K$-a.e. for every $\\rho \\in\\mathcal P$.\n\n{\\it Step 3.}In view of {\\it Steps 1 and 2}, uniform integrability of $\\mathcal P$ under $dK$ follows from the integrability of $f$ under $d K$.\n\\end{proof}\n\n\\begin{Lemma}\\label{uiB}\nLet the assumption of Theorem \\ref{mainTheorem2} hold and $(x,q)\\in\\mathring{\\mathcal K}$. Then \n$$\\left\\{r:(y,r)\\in\\mathcal B(x,q)\\right\\}$$ is \n $\\mathbb{L}^{\\infty}$-bounded, so a uniformly integrable subset of \n$\\mathbb L^1(dK)$.\n\\end{Lemma}\n\\begin{proof}\nBy Lemma \\ref{7-17-1}, we deduce the existence of a constant $M>0$,\nsuch that\n\\begin{equation}\\nonumbe\ny\\leq M,\\quad for~every~(y,r)\\in\\mathcal B(x,q).\n\\end{equation}\nWe conclude that $$\\left\\{r:(y,r)\\in\\mathcal B(x,q)\\right\\}\\subseteq \\bigcup\\limits_{0\\leq \\lambda \\leq M}{\\lambda\\mathcal P},$$ and thus by Lemma \\ref{uiP}, $\\left\\{r:(y,r)\\in\\mathcal B(x,q)\\right\\}$ is a uniformly integrable family.\n\\end{proof}\n\\subsection{Closedness of $\\mathcal D(x,q)$ for every $(x,q)\\in\\mathring{\\mathcal K}$}\n\\begin{Lemma}\\label{closednessD}\nUnder the conditions of Theorem~\\ref{mainTheorem2}, for every $(x,q)\\in \\mathring{\\mathcal K}$, the set $\\mathcal D(x,q)$ is closed in $\\mathbb L^0(d\\kappa\\times\\mathbb P)$.\n\\end{Lemma}\n\\begin{proof}\nLet $(x,q)\\in\\mathring{\\mathcal K}$\nand $Y$ be an arbitrary element of $\\cl \\mathcal D(x,q)$. We claim that there exists $(y,r)\\in \\mathcal B(x,q)$, such that $Y \\in \\mathcal Y(y,r)$. Let $Y^n\\in\\mathcal Y(y^n,r^n)$, $n\\geq 1$, be a sequence in $\\mathcal D(x,q)$, such that $\\lim\\limits_{n\\to \\infty}Y^n = Y$, $(d\\kappa\\times\\mathbb P)$-a.e. \nSince $(y^n,r^n)_{n\\geq 1}\\subset\\mathcal B(x,q)$, which is bounded in the sense of Lemma~\\ref{7-17-1}, Komlos' lemma implies the existence of a subsequence of convex combinations\n$(\\tilde y^n,\\tilde r^n)\\in \\conv((y^n,r^n), (y^{n+1}, r^{n+1}),\\dots)$, $n\\geq 1$, such that $(\\tilde y^n)_{n\\geq 1}$ converges to $y$ and $(\\tilde r^n)_{n\\geq 1}$ converges to $r$, $dK$-a.e. \nLemma~\\ref{uiB} implies that $(\\tilde r^n)_{n\\geq 1}$ is uniformly integrable. Therefore $(\\tilde r^n)_{n\\geq 1}$ converges to $r$ in $\\mathbb L^1(dK)$. Note that the corresponding sequence of convex combinations of $(Y^n)_{n\\geq 1}$, $(\\tilde Y^n)_{n\\geq 1}$ converges to $Y$, $(d\\kappa\\times\\mathbb P)$-a.e.\nThen we have\n$$ \n1 \\geq \\lim\\limits_{n\\to\\infty}\\left( x\\tilde y^n + \\int_0^T q_s\\tilde r^n_sdK_s \\right) = xy + \\int_0^T q_s r_sdK_s.\n$$\nTherefore, $(y,r)\\in \\mathcal B(x,q)$.\n\nLet us fix an arbitrary $(x',q')\\in\\mathcal K$ and $c\\in\\mathcal A(x',q')$. \nUsing Proposition~\\ref{prop1}, Fatou's lemma, and Lemma \\ref{uiB}, we obtain\n\\begin{displaymath}\n\\begin{array}{c}\n0\\leq \\mathbb E\\left[\\int_0^TY_sc_sds \\right] \\leq \n\\liminf\\limits_{n\\to\\infty} \\mathbb E\\left[\\int_0^T \\tilde Y^n_sc_sds \\right] \\\\\\leq \n \\lim\\limits_{n\\to\\infty}\\left( x'\\tilde y^n + \\int_0^T q'_s\\tilde r^n_sdK_s \\right)\n =\n x'y + \\int_0^T q'_s r_sdK_s,\\\\\n\\end{array}\n\\end{displaymath} \nwhere the uniform integrability of $\\mathcal B(x,q)$ in needed once again in the last equality.\nFrom Proposition~\\ref{prop1}, we conclude that that $Y\\in \\mathcal Y(y,r)$. \n\\end{proof}\n\n\\begin{proof}[Proof of Theorem~\\ref{mainTheorem2}.] \nLet $(x,q)\\in\\mathcal K$, $\\widehat c(x,q)$ be the minimizer to (\\ref{primalProblem}), whose existence and uniqueness are established in Lemma~\\ref{existenceUniqueness}. Let us also set\n\\begin{equation}\\label{defHatY}\n\\widehat Y_{t} {:=} U'({t},\\widehat c_{t}(x,q)),\\quad (t,\\omega)\\in[0,T]\\times\\Omega,\\quad and \\quad z {:=} \\mathbb E\\left[\\int_0^T\\widehat c_s(x,q)\\widehat Y_s d\\kappa_s \\right].\n\\end{equation} \nNote that the sets $\\mathcal A(x,q)$ and ${\\cl \\mathcal D(x,q)}$ satisfy the conditions of \\cite[Theorem~3.2]{MostovyiNec}, which implies that $\\widehat Y\\in z{\\cl \\mathcal D(x,q)}$ is the unique solution to the optimization problem\n\\begin{displaymath}\n\\inf\\limits_{Y\\in z{\\cl \\mathcal D(x,q)}}\\mathbb E\\left[\\int_0^TV(t, Y_s)d\\kappa_s\\right] = \\mathbb E\\left[\\int_0^TV(t,\\widehat Y_s)d\\kappa_s\\right]\\in\\mathbb R,\n\\end{displaymath}\nwhere finiteness follows from Lemma \\ref{auxiliaryLemma}.\nNote that by Lemma~\\ref{closednessD}, $\\widehat Y\\in\\mathcal Y(zy, zr)$ for some $(y,r)\\in\\mathcal B(x,q)$.\nIt follows from the definition of $\\widehat Y$ in \\eqref{defHatY} that \nfor $\\widehat c(x,q)$ and $\\widehat Y$ we have the following relation\n\\begin{displaymath}\nU({t}, \\widehat c_{t}(x,q)) = V({t}, \\widehat Y_{t}) + \\widehat c_{t}(x,q)\\widehat Y_{t},\n\\quad (t,\\omega)\\in[0,T]\\times\\Omega\n\\end{displaymath}\nwhich together with Lemma~\\ref{conjugacy} implies that\n\\begin{equation}\\label{7-17-3}\nu(x,q) = v(zy,zr) + z\\left(xy + \\int_0^Tq_sr_sdK_s\\right),\n\\end{equation}\n\\begin{displaymath}\n\\widehat Y(zy,zr) = \\widehat Y\n,\\quad (d\\kappa\\times\\mathbb P)-a.e.,\n\\end{displaymath}\nwhere $\\widehat Y(zy,zr)$ is the unique minimizer to the dual problem~(\\ref{dualProblem}). \nBy (\\ref{7-17-3}), $(zy,zr)\\in\\mathcal E.$ (\\ref{7-17-3}) and \\cite[Proposition I.5.1]{EkelandTemam} assert that $(zy,zr)\\in\\partial u(x,q)$. In particular, we get\n\\begin{displaymath} \n\\partial u(x,q) \\cap \\mathcal E \\neq \\emptyset,\n\\end{displaymath}\ni.e., \\eqref{subdifNonempty}. \nNote that even though, e.g., \\cite[Corollary 2.2.38 and Corollary 2.2.44]{BarbuPrec} imply that $\\partial u(x,q)\\neq \\emptyset$, \n over the interior of the effective domain, \nits elements are in $\\mathbb R\\times (\\mathbb L^\\infty)^*(dK)$. Relation \\eqref{subdifNonempty} shows that $\\partial u(x,q)$ contains at least a bounded element of $\\mathcal E\\subseteq\\mathcal L\\subset \\mathbb R\\times\\mathbb L^1(dK).$\n\nLet $(x,q)\\in\\mathcal K$ and $(y,r)\\in\\mathcal L$. Suppose that \\eqref{2191}, \\eqref{2192}, and \\eqref{2193} hold. Then by conjugacy of $U$ and $V$, we get \n\n\\begin{equation}\\label{2194}\n0 = v(y,r) - u(x,q) + xy + \\int_0^T q_sr_sdK_s.\n\\end{equation}\nLemma \\ref{lemBiconjugacy} and \n\\cite[Proposition I.5.1]{EkelandTemam} imply that $(y,r)\\in\\partial u(x,q) \\cap\\mathcal E$.\nConversely, let $(x,q)\\in\\mathcal K$ and $(y,r)\\in\\mathcal L\\cap \\partial u(x,q).$ Then by \\cite[Proposition I.5.1]{EkelandTemam} and Lemma \\ref{lemBiconjugacy}, we deduce that \\eqref{2194} holds. Lemma \\ref{lemFinu} implies the finiteness of $u(x,q)$, which together with \\eqref{2194} results in the finiteness of $v(y,r)$, thus \\eqref{2191} holds and $(y,r)\\in\\mathcal E$. By Lemma \\ref{existenceUniqueness}, there exists a unique optimizer $\\widehat c(x,q)$, for \\eqref{primalProblem}, and $\\widehat Y(y,r)$, for \\eqref{dualProblem}, respectively. Therefore, from conjugacy of $U$ and $V$, Proposition \\ref{prop1}, and \\eqref{2194}, we obtain\n\\begin{displaymath}\n\\begin{array}{c}\n0\\leq\n \\mathbb E\\left[\\int_0^T \\left(V(t,\\widehat Y_t(y,r)) - U(t,\\widehat c_t(x,q)) + \\widehat Y_t(y,r) \\widehat c_t(x,q)\\right)d\\kappa_t\\right] \\\\\n\\leq v(y,r) - u(x,q) + xy + \\int_0^T q_sr_sdK_s = 0.\\\\\n\\end{array}\n\\end{displaymath}\nThis implies \\eqref{2192} and \\eqref{2193}. This completes the proof of the theorem.\n\\end{proof}\n\n\\section{Structure of the dual feasible set}\\label{secStructureOfDualDomain}\nBy Assumption \\ref{asACclock}, there exists at most countable subset $(s_k)_{k\\in\\mathbb N}$ of $[0,T]$, where $\\kappa$ has jumps. \nWe define $\\mathcal{D}'$ the set of non-increasing, left-continuous and adapted processes $D$ that start at $1$ and \nwith the property that \n$D_{\\theta _0} = 1$, $ D_T\\geq 0$\nand that, there exists some $n\\in \\mathbb N$ such that $D$ is constant off the discrete grid\n$\\mathcal T_n{:=} \\bigcup\\limits_{j = 1}^{2^n}\\{s_j\\}\\cup\\left\\{\\frac{k}{2^n}T, k = 0, \\dots,2^n\\right\\}$.\n\n\n\n\n\n\\begin{Lemma}\\label{lemSt2}\n Let $\\mathcal G'{:=} \\{ZD = (Z_tD_t)_{t\\in[0,T]}:~Z\\in\\mathcal Z', D\\in\\mathcal D' \\}$. Assume the conditions of Proposition \\ref{prop1} hold. Then $\\mathcal G'$ is convex.\n\n\\end{Lemma}\n\\begin{proof}\nLet $Z^1D^1$ and $Z^2D^2$ are the elements of $\\mathcal G'$ and let $\\lambda \\in(0,1)$\nWe need to show that $\\lambda Z^1D^1 + (1-\\lambda)Z^2D^2 = ZD$\nfor some $Z\\in\\mathcal Z'$ and $D\\in\\mathcal D'$. There exists $n\\in\\mathbb N$, such that $D^1$ and $D^2$ decrease at most on $\\mathcal T_n$. Let $t_k$'s be the elements of $\\mathcal T_n$ arranged in an increasing order. \nLet us define $Z_0 = D_0 = 1$ and for every $k \\in\\{0,\\dots, 2^n-1\\}$,~with \n$$\nA_t {:=} \\lambda Z^1_{t_k}D^1_{t} +(1-\\lambda)Z^2_{t_k}D^2_{t}\n\\quad and\\quad\n\\alpha_t {:=} \\frac{\\lambda Z^1_{t_k}D^1_{t}}{A_t}1_{\\{A_t\\neq 0\\}} + \\frac{\\lambda Z^1_{t_k}}{\\lambda Z^1_{t_k} + (1-\\lambda) Z^2_{t_k}}1_{\\{A_t= 0\\}},$$\n(note that $A_t = 0$ if and only if both $D^1_t=0$ and $D^2_t=0$) we set\n\\begin{equation}\\label{351}\n\\begin{array}{rclcl}\nZ_t &{:=}& Z_{t_k}\\left(\\alpha_{t_k+}\\frac{Z^1_t}{Z^1_{t_k}} + (1 - \\alpha_{t_k+})\\frac{Z^2_t}{Z^2_{t_k}} \\right), &for& t\\in(t_k, t_{k+1}],\\\\\nD_t &{:=}& \\frac{A_{t_k+}}{Z_{t_k}},& for& t\\in(t_k, t_{k+1}].\\\\\n\\end{array}\n\\end{equation}\nOne can see that $ZD = \\lambda Z^1D^1 + (1-\\lambda)Z^2D^2$ and that $Z\\in\\mathcal Z'$, see e.g., \\cite{FK}, \\cite{Rohlin10}, \\cite{Kardaras_emery}, and \\cite{ChristaJosef15} for discussions of the sets of processes with similar convexity-type properties to the one given in \\eqref{351}. To show that $D\\in\\mathcal D'$, for $k\\geq 1$, \nwe observe that \n\\begin{displaymath}\n\\begin{array}{rcl}\nD_{t_k+} &=& \\frac{A_{t_k+}}{Z_{t_k}} \\\\\n&=& \\frac{{\\lambda Z^1_{t_k}D^1_{t_k+} +(1-\\lambda)Z^2_{t_k}D^2_{t_k+}}}{Z_{t_{k-1}}\\left(\\alpha_{t_{k-1}+}\\frac{Z^1_{t_k}}{Z^1_{t_{k-1}}} + (1 - \\alpha_{t_{k-1}+})\\frac{Z^2_{t_k}}{Z^2_{t_{k-1}}} \\right)} \\\\\n&\\leq &\\frac{{\\lambda Z^1_{t_k}D^1_{t_k} +(1-\\lambda)Z^2_{t_k}D^2_{t_k}}}{Z_{t_{k-1}}\\left(\\alpha_{t_{k-1}+}\\frac{Z^1_{t_k}}{Z^1_{t_{k-1}}} + (1 - \\alpha_{t_{k-1}+})\\frac{Z^2_{t_k}}{Z^2_{t_{k-1}}} \\right)}\\\\\n&= &\\frac{{\\lambda Z^1_{t_k}D^1_{t_k} +(1-\\lambda)Z^2_{t_k}D^2_{t_k}}}{Z_{t_{k-1}}\\left(\\frac{\\lambda Z^1_{t_{k-1}}D^1_{t_k}}{A_{t_{k-1}+}}\\frac{Z^1_{t_k}}{Z^1_{t_{k-1}}} \n+ \\frac{(1-\\lambda) Z^2_{t_{k-1}}D^2_{t_k}}{A_{t_{k-1}+}}\\frac{Z^2_{t_k}}{Z^2_{t_{k-1}}} \\right)}1_{\\{A_{t_{k-1}+}> 0\\}} +0\\cdot1_{\\{A_{t_{k-1}+}= 0\\}}\\\\\n&= &\\frac{A_{t_{k-1}+}}{Z_{t_{k-1}}}\\frac{{\\lambda Z^1_{t_k}D^1_{t_k} +(1-\\lambda)Z^2_{t_k}D^2_{t_k}}}{\\lambda D^1_{t_k}Z^1_{t_k}\n+ (1-\\lambda) D^2_{t_k}Z^2_{t_k}}1_{\\{A_{t_{k-1}+}> 0\\}}\\\\\n&= &\\frac{A_{t_{k-1}+}}{Z_{t_{k-1}}}1_{\\{A_{t_{k-1}+}> 0\\}}\\\\\n&= &D_{t_k}.\\\\\n\\end{array}\n\\end{displaymath}\nAbove, the inequality follows from the monotonicity of $D^1$ and $D^2$.\nTherefore, $D$ is nonincreasing. Also, clearly $D$ is nonnegative. Thus $ZD\\in \\mathcal G'$. \n\\end{proof}\nThe following lemma is an extension of Lemma \\ref{9-5-1} and amounts to a first layer of convexification of the set $\\Upsilon$, i.e., of the budget constraints.\n\n\\begin{Lemma}\\label{lemSt1}Let the conditions of Proposition \\ref{prop1} hold, $(x,q)\\in\\mathcal K$, and $c$ is a nonnegative optional process. Then $c\\in\\mathcal A(x,q)$ if and only if\n\\begin{equation}\\label{331}\n\\mathbb E\\left[\\int_0^Tc_sD_sZ_sd\\kappa_s \\right]\\leq\nx + \\mathbb E\\left[\\int_0^Tq_se_sD_sZ_sd\\kappa_s \\right],\\quad\nfor~every~Z\\in\\mathcal Z'~and~D \\in \\mathcal D'.\n\\end{equation}\n\\end{Lemma}\n\\begin{proof}\nThe idea is to use the assertion of Lemma \\ref{9-5-1} and to approximate a given $D\\in\\mathcal D'$ by (finite) linear combinations of the elements of $\\Upsilon$, where $\\Upsilon$ is defined in \\eqref{defUpsilon}.\n\nFor a stopping time $\\tau$, \nlet us denote $$\\Lambda^{\\tau} {:=} 1_{[0,\\tau]}\\in\\Upsilon,$$\nand fix $D\\in\\mathcal D'$. Then there exists $l\\in\\mathbb N$, such that $D$ has has jumps at most\non $\\{t_0,t_1,\\dots,\nt_l\\}$ for some increasing $t_i$'s. \nFor every $j \\in\\{ 0, \\dots, \n\\}$, $k \\in\\{0,\\dots 2^n\\}$, and $n\\in\\mathbb N$, let us set\n\\begin{equation}\\nonumbe\n\\begin{array}{rcl}\nA_{k, n, j} &{:=}& \\left\\{ \\omega: ~D_{t_j}(\\omega) >0~and~\\frac{D_{t_j+}(\\omega)}{D_{t_j}(\\omega)} \\in \\left(\\frac{k-1}{2^n}, \\frac{k}{2^n}\\right]\\right\\},\\\\\n\\tau^{k, n, j}&{:=}&T1_{A_{k,n,j}} + t_j1_{A^c_{k,n,j}},\\\\\n\\end{array}\n\\end{equation}\nNote that $D_0 = 1$ by definition of $\\mathcal D'$, $A_{k, n, j}\\in\\mathcal F_{t_j}$, and \n$$\\frac{k-1}{2^n}\\Lambda^{\\tau^{k,n,j}}_{t_j+} = \\left\\{\n\\begin{array}{lcl}\n\\frac{k-1}{2^n}&on&A_{k,n,j}\\\\\n0&on& A^c_{k,n,j}\\\\\n\\end{array}\n \\right..$$\n By construction, we have $$D_{t_j+} = D_{t_j}\\lim\\limits_{n\\to\\infty}\\sum\\limits_{k = 1}^{2^n}\\frac{k-1}{2^n}\\Lambda^{\\tau^{k,n,j}}_{t_j+},$$ where the sequence \n $$K^n_{t_j+}{:=} \\sum\\limits_{k = 1}^{2^n}\\frac{k-1}{2^n}\\Lambda^{k,n,j}_{t_j+},\\quad n\\in\\mathbb N,$$\n is increasing on $\\{D_{t_j}>0\\}$, i.e.,\n $$K^n_{t_j+}1_{\\{D_{t_j}>0\\}}\\uparrow \\frac{D_{t_j+}}{D_{t_j}}1_{\\{D_{t_j}>0\\}}.$$\n Thus, for an arbitrary $j \\in\\{0, \\dots, l\\}$, we have constructed a sequence of elements of $\\Upsilon$, whose finite linear combinations monotonically increase to $\\frac{D_{t_j+}}{D_{t_j}}$ on $\\{D_{t_j}>0\\}$ (i.e., if $j0\\}$, we have\n$$\\lim\\limits_{n\\to\\infty}\\sum\\limits_{i=1}^{4^n}\\lambda^{n,i} \\Lambda^{\\sigma_{n,i}}_{t+} = \\frac{D_{t+}}{D_t}.$$\nSimilarly,\n with $r(t){:=} \\max\\{i:~t_i< t\\}$, let us define\n$$D^n_t {:=} \\prod\\limits_{j = 0}^{r(t)}K^n_{t_j+}1_{\\{D_{t_j}>0\\}},\\quad t\\in[0,T], n\\in\\mathbb N.$$ \nAs in \\eqref{361}, for every $n\\in\\mathbb N$, $D^n$ can be written as a finite linear combination of $\\Lambda$'s, such that $D^n_{t+}\\uparrow D_{t+}$ for every $t\\in\\{t_0, \\dots, t_{l}\\}.$\n\nFinally, \\eqref{331} can be obtained from Lemma \\ref{9-5-1} by the approximation of $D$ by $D^n$'s as above and via the monotone convergence theorem (applied separately to $(e_t)^{+}$ and $(e_t)^{-}$). \n\\end{proof}\n\\begin{Corollary}\\label{corRzd}\nLet the conditions of Proposition \\ref{prop1} hold. Then, for every pair $Z\\in\\mathcal Z'$ and $D\\in\\mathcal D'$, \nthere exists $r^{ZD}\\in\\mathbb L^1(dK)$, such that $ZD\\in\\mathcal Y(1,r^{ZD})$, where $(1,r^{ZD})\\in\\mathcal L$. \n\\end{Corollary}\n\\begin{proof}\nThe existence of $r^{ZD}$, such that $ZD\\in\\mathcal Y(1, r^{ZD})$ follows from Lemma \\ref{lemSt1} (equation \\eqref{331}) and the approximation procedure in Lemma \\ref{lemSt1} (again, applied separately to $(e_t)^{+}$ and $(e_t)^{-}$) combined with the monotone convergence theorem.\nAs, the left-hand side in \\eqref{331} is nonnegative, \n$(1,r^{ZD})\\in\\mathcal L$. \n\n\\end{proof}\n\n\nFor a given $Z\\in\\mathcal Z'$ and $D\\in\\mathcal D'$, let us recall that $r^{ZD}$ is given in Corollary \\ref{corRzd}. \nWe set $$\\mathcal {B'}(x,q) {:=} \\left\\{ (y, yr^{ZD})\\in\\mathcal B(x,q):~y> 0, Z\\in\\mathcal Z', D\\in\\mathcal D'\\right\\},$$\n\n$$\\mathcal G(x,q) {:=} \\left\\{yZ'D'\\in\\mathcal D(x,q):~(y, yr^{ZD})\\in \\mathcal {B'}(x,q)\\right\\},\\quad (x,q)\\in\\mathring{\\mathcal K},$$ \n\n\n\\begin{Lemma}\\label{lemSt3}\nLet the conditions of Proposition \\ref{prop1} hold, then for every $(x,q) \\in\\mathring{\\mathcal K}$, the closure of the convex, solid hull of $\\mathcal G(x,q)$ in $\\mathbb L^0$ coincides with $\\cl\\mathcal D(x,q)$. \n\\end{Lemma}\n\\begin{proof}\nLet $(x,q)\\in\\mathring{\\mathcal K}$ be fixed. \nAlong the lines of the proof of Lemma \\ref{5181}, one can show that\n\\begin{equation}\\nonumbe\n(\\mathcal G(x,q))^o = \\mathcal A(x,q).\n\\end{equation}\nTherefore, \n$$(\\mathcal G(x,q))^{oo} = (\\mathcal A(x,q))^o = \\cl\\mathcal D(x,q),$$\nwhere in the last equality we have used the conclusion of Lemma \\ref{5181}.\n As $\\mathcal G(x,q)\\subset \\cl\\mathcal D(x,q)$, the assertion of the lemma follows from the bipolar theorem of Brannath and Schchermayer, \\cite[Theorem 1.3]{BranSchach}.\n\\end{proof}\n\n\n\\begin{Corollary}\\label{lemSt4}\nLet the conditions of Proposition \\ref{prop1} hold and $(x,q)\\in\\mathring{\\mathcal K}$. Then for \nevery \\underline{maximal} element $Y$ of $\\cl\\mathcal D(x,q)$ there exists $y^n\\geq 0$, $Z^n\\in\\mathcal Z'$, $D^n\\in\\mathcal D'$, $n\\in\\mathbb N$, such that $(y^nZ^nD^n)_{n\\in\\mathbb N}\\subset\\mathcal G(x,q)$ and \n\\begin{equation}\\nonumber\n\\begin{array}{rcll}\nY\n&=& \\lim\\limits_{n\\to\\infty} y^nZ^nD^n,& \\quad (d\\kappa\\times \\mathbb P)-a.e.~and~on~ \\bigcup\\limits_{n\\in\\mathbb N}\\mathcal T_n.\n\\end{array}\n\\end{equation}\n\\end{Corollary}\n\\begin{proof}\nThe $(d\\kappa\\times \\mathbb P)$-a.e. convergence follows from Lemma \\ref{lemSt3}. By passing to subsequences of convex combinations, we also deduce the convergence on $\\bigcup\\limits_{n\\in\\mathbb N}\\mathcal T_n$.\n\\end{proof}\n\\begin{Remark}\nIt follows from Corollary \\ref{lemSt4} and Fatou's lemma that the maximal elements of $\\cl\\mathcal D(x,q)$ are strong supermartingales. Moreover, every maximal element of $\\cl\\mathcal D(x,q)$ is an optional strong supermartingale deflator, which is optional strong supermartingale $Y$, such that $XY$ is an optional strong supermartingale for every $X\\in\\mathcal X(1)$. We refer to \\cite[Appendix 1]{DelMey82} for a general characterization and to \\cite{CzichSchachOSSup} for results on strong optional supermartingales as limits of martingales. The following section gives a more refined characterization of the dual minimizer.\n\\end{Remark}\n\n\n\\section{Complementary slackness}\\label{secSlackness}\nFor better readability of this section, we recall some notations and results that will be used below. \nThroughout this section, $(x,q)\\in \\mathring{\\mathcal{K}}$ will be fixed, $\\widehat c = \\widehat c(x,q)$ is the optimizer to \\eqref{primalProblem}, $\\widehat V$ is the corresponding wealth process, i.e., \n \\begin{equation}\\label{5211}\n\\widehat V=x+\\int_0^{\\cdot}\\widehat H_sdS_s-\\int_0^{\\cdot}\\widehat c_sd\\kappa_s+\\int_0^{\\cdot} q_se_sd\\kappa_s,\n\\end{equation}\nwhere $\\widehat H$ is some $S$-integrable process, \n$\\widehat Y$ be such that $\\widehat Y_t = U'(t,\\widehat c_t)$, $(d\\kappa\\times\\mathbb P)$-a.e., i.e., $\\widehat Y$ is the optimizer to \\eqref{dualProblem} for some \n$(y,r)\\in\\mathcal E\\cap \\partial u(x,q)$, $\\widehat Y\\in\\mathcal Y(y, r)$ and \n$$\\mathbb E\\left[\\int_0^TV(t,\\widehat Y_t)d\\kappa_t\\right] = \\inf\\limits_{Y\\in \\mathcal Y(y,r)}\\mathbb E\\left[\\int_0^TV(t, Y_t)d\\kappa_t\\right].$$\nBy Corollary \\ref{lemSt4}, $\\widehat Y$ can be approximated by a sequence $y^nD^nZ^n$, $n\\in\\mathbb N$, where $y^n$ is a nonnegative constant, $D^n\\in\\mathcal D'$, and $Z^n\\in\\mathcal Z'$, $n\\in\\mathbb N$.\nIn what follows, for any right-continuous increasing process $A$ satisfying $$A_t=0,\\quad 0-\\leq t\\leq \\theta _{0}-,\\quad A_T=1,$$ \ni.e., for any probability measure $dA$ on the closed interval $[\\theta _0, T]$ (extended to $[0,T]$) we will associate a process $D$ which is left-continuous and decreasing\n$$D_t{:=} 1-A_{t-}=dA([t,T]), \\ \\ 0\\leq t\\leq T.$$\nOne can also think that $ D_{T+}=0,$ although this is not necessary. It is clear that such $A\\leftrightarrow D$ are in bijective correspondence. Below, all processes $A$'s and $D$'s (with indexes) will be in such bijective correspondence, except for the case of the limiting process $\\widehat A$ (which is right-continuous) and the limiting process $\\widehat D$ (that may be not left-continuous). The will be in a similar but more subtle correspondence. More precisely:\n\\begin{Theorem}\\label{ExactSlackness}\nLet the conditions of Theorem \\ref{mainTheorem2} hold. Let $\\widehat V$ be the optimal wealth process, $\\widehat c$ the optimal consumption, and $\\widehat Y$ be the dual minimizer satisfying the assertions of Theorem \\ref{mainTheorem2}. Then, there exists a strong supermartingale $\\hat Z=\\lim\\limits_{n\\to\\infty} Z^n$ (for $Z_n\\in \\mathcal{D}'$, where the limit is in the sense of \\cite{CzichSchachOSSup}) and a right-continuous increasing process $\\widehat A$ with\n$$ \\widehat A_t=0\\ \\ \\forall 0-\\leq t \\leq \\theta _0-,\\ \\ \\ \\widehat A_T= 1, $$\nand a decreasing process $\\widehat D$ with $\\widehat D_t=1,\\ \\ 0\\leq t\\leq {\\theta _0}$ and satisfying the complementary slackness condition \n\\begin{equation}\\label{comp-slack}\n \\mathbb{P}\\left (\\widehat D_t\\in \\left[1-\\widehat A_t, 1-\\widehat A_{t-}\\right], \\ \\forall \\ \\theta _0 \\leq t\\leq T)\\right)=1,\\ \\ \n \\mathbb{P}\\left (\\int _ {[\\theta _0,T)}1_{\\{\\widehat V_{t-}\\not= 0, \\widehat V_t \\not =0\\}} d\\widehat A_t\\right)=0,\n \\end{equation}\nsuch that the dual minimizer can be decomposed as\n$$\\hat Y= y \\hat Z\\hat D.$$\n\\end{Theorem}\n\nThe proof of the Theorem \\ref{ExactSlackness} is split in several results. \n\\begin{Lemma}\\label{lem5211} Let the conditions of Theorem \\ref{mainTheorem2} hold. With \n$$z{:=} xy+\\int_0^T q_sr_sdK_s,$$ there exist $y^n>0$, $Z^n\\in\\mathcal Z'$, and $D^n\\in\\mathcal D'$, such that $y^nZ^nD^n\\in z\\mathcal G(x,q)$, $n\\in\\mathbb N$, and\n\\begin{equation}\\label{5231}\n\\begin{array}{c}\n\\widehat Y\n= \\lim\\limits_{n\\to\\infty} y^nZ^nD^n,\\quad (d\\kappa\\times \\mathbb P)-a.e.\\\\\ny^n\\to y>0;\\\\%\\textcolor{blue}{\\quad Z^n\\to \\widehat Z,\\quad D^n\\to \\widehat D,\\quad (d\\kappa\\times \\mathbb P)-a.e.}\\\\\n\\end{array}\n\\end{equation}\nFor $A^n$, $n\\in\\mathbb N$, being in relation to $D^n$ exactly as described before Theorem \\ref{ExactSlackness} we have\n\\begin{equation}\\label{86*}\\mathbb{E}\\left [\\int _{\\theta _0}^T \\widehat V_t Z^n_t dA^n _t \\right]=\\mathbb{E}\\left [Z^n_T\\int _{\\theta _0}^T \\widehat V_t dA^n _t \\right]\\rightarrow 0.\n\\end{equation}\n\\end{Lemma}\n\\begin{proof} Optimality of $\\widehat Y$ and Corollary \\ref{lemSt4} imply \\eqref{5231}. \nBy \\eqref{2192} and Fatou's lemma, we get\n\\begin{equation}\\label{eqndefz}\nz = \\mathbb E\\left[\\int_0^T \\widehat Y_s\\widehat c_sd\\kappa_s \\right] \\leq \\liminf\\limits_{n\\to\\infty}y^n\\mathbb E\\left[ \\int_0^T Z^n_sD^n_s\\widehat c_sd\\kappa_s\\right].\n\\end{equation}\nLet us fix $n\\in\\mathbb N$ and consider $\\mathbb E\\left[ \\int_0^T Z^n_sD^n_s\\widehat c_sd\\kappa_s\\right]$. Using localization and integration by parts (along the lines of the proof of Lemma \\ref{6-16-1}), we have\n\\begin{displaymath}\\begin{array}{rcl}\n\\mathbb E\\left[ \\int_0^T Z^n_sD^n_s\\widehat c_sd\\kappa_s\\right] &=& \\mathbb E\\left[Z^n_T \\int_0^T D^n_{s}\\widehat c_sd\\kappa_s\\right] \\\\\n&=& \\mathbb E\\left[Z^n_T \\int_0^T \\left(1 - A^n_{s-}\\right)\\widehat c_sd\\kappa_s\\right] \\\\\n&=& \\mathbb E\\left[Z^n_T \\left( \\int_0^T\\widehat c_sd\\kappa_s - \\int_0^T A^n_{s-} \\widehat c_sd\\kappa_s\\right)\\right] \\\\\n&=& \\mathbb E\\left[Z^n_T \\left( \\int_0^T\\widehat c_sd\\kappa_s - A^n_{T} \\int_0^T \\widehat c_sd\\kappa_s + \\int_0^T(\\int_0^t\\widehat c_sd\\kappa_s)dA^n_t \\right)\\right] \\\\\n&=& \\mathbb E\\left[Z^n_T \\int_0^T(\\int_0^t\\widehat c_sd\\kappa_s)dA^n_t \\right].\\\\\n\\end{array}\n\\end{displaymath}\nThe latter expression, using \\eqref{5211} and with $\\bar X {:=} ||q||_{\\mathbb L^\\infty(dK)}X''$ (where in turn $X''$ is given by the assertion $(v)$ of Lemma \\ref{lemma1}), we can rewrite as\n\\begin{equation}\\label{5212}\n\\mathbb E\\left[Z^n_T \\int_0^T\\left((x + \\int_0^{t}\\widehat H_sdS_s +\\bar X_t) + \\left(\\int_0^tq_se_sd\\kappa_s - \\bar X_t\\right) - \\widehat V_t\\right)dA^n_t \\right].\n\\end{equation}\nLet us denote\n\\begin{equation}\\nonumbe\n\\begin{array}\n{rcl}\nT_1 &{:=} & \\mathbb E\\left[Z^n_T \\int_0^T\\left(x + \\int_0^{t}\\widehat H_sdS_s +\\bar X_t\\right)dA^n_t \\right] ,\\\\\nT_2 &{:=} & \\mathbb E\\left[Z^n_T \\int_0^T\\left(\\int_0^tq_se_sd\\kappa_s - \\bar X_t\\right)dA^n_t \\right].\\\\\n\\end{array}\n\\end{equation}\nIt follows from nonnegativity of $\\widehat V$ in \\eqref{5211}, Lemma \\ref{lemma1}, and nonnegativity of $\\widehat c$ that \n\\begin{equation}\\label{5232}\nx + \\int_0^{t}\\widehat H_sdS_s + \\bar X_t \\geq \\int_0^t\\widehat c_sd\\kappa_s + \\bar X_t -\\int_0^t q_se_sd\\kappa_s \\geq \\bar X_t -\\int_0^t q_se_sd\\kappa_s \\geq0,\\quad t\\in[0,T],\n\\end{equation}\ni.e., the integrand in $T_1$ is nonnegative. Let ${\\mathbb Q^n}$ be the probability measure, whose density process with respect to $\\mathbb P$ is $Z^n$. As $(x + \\int_0^{\\cdot}\\widehat H_sdS_s + \\bar X)$ and $\\bar X$ are local martingales under ${\\mathbb Q^n}$, by \\cite[Theorem III.27, p. 128]{Pr}, $A^n_{-}\\cdot(x + \\int_0^{\\cdot}\\widehat H_sdS_s + \\bar X)$ and $A^n_{-}\\cdot\\bar X$ are local martingales. Let $(\\tau_k)_{k\\in\\mathbb N}$ be a localizing sequence for both $A^n_{-}\\cdot(x + \\int_0^{\\cdot}\\widehat H_sdS_s + \\bar X)$ and $A^n_{-}\\cdot\\bar X$. \nBy the monotone convergence theorem and integration by parts, we get\n\\begin{equation}\\label{5214}\n\\begin{array}{rcl}\nT_1 &=& \\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[\\int_0^{\\tau_k}\\left(x + \\int_0^{t}\\widehat H_sdS_s +\\bar X_t\\right)dA^n_t \\right] \\\\\n&=& \\lim\\limits_{k\\to\\infty}\\left(\\mathbb E^{\\mathbb Q^n}\\left[ \\int_0^{\\tau_k}(-A^n_{t-})d(x + \\int_0^{t}\\widehat H_sdS_s + \\bar X_t) + A^n_{\\tau_k} \\left( x + \\int_0^{\\tau_k}\\widehat H_sdS_s +\\bar X_{\\tau_k}\\right)\\right]\\right) \\\\\n&=&\\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[A^n_{\\tau_k} \\left( x + \\int_0^{\\tau_k}\\widehat H_sdS_s +\\bar X_{\\tau_k}\\right)\\right].\\\\\n\\end{array}\n\\end{equation}\nLet us consider $T_2$. With $\\mathcal E^q {:=} \\int_0^{\\cdot} q_se_sd\\kappa_s$, Lemma \\ref{lemma1} implies positivity of $\\mathcal E^q_t - \\bar X_t$, $t\\in[0,T]$, which in turn allows to invoke the monotone convergence theorem, and we obtain \n\\begin{equation}\\nonumbe\n\\begin{array}{rcl}\nT_2 &=& \\mathbb E^{\\mathbb Q^n}\\left[ \\int_0^T\\left(\\mathcal E^q_t - \\bar X_t\\right)dA^n_t \\right] \\\\\n&=&\\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[ \\int_0^{\\tau_k}\\mathcal E^q_t dA^n_t - \\int_0^{\\tau_k}\\bar X_tdA^n_t \\right].\\\\\n\\end{array}\n\\end{equation} \nBy positivity of $\\bar X$ and the monotone convergence theorem, we have $$\\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[\\int_0^{\\tau_k}\\bar X_tdA^n_t \\right]=\\mathbb E^{\\mathbb Q^n}\\left[\\int_0^{T}\\bar X_tdA^n_t \\right].$$\nTherefore, $\\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[ \\int_0^{\\tau_k}\\mathcal E^q_t dA^n_t\\right]=\\mathbb E^{\\mathbb Q^n}\\left[ \\int_0^{T}\\mathcal E^q_t dA^n_t\\right]$. We deduce that \n\\begin{displaymath}\n\\mathbb E^{\\mathbb Q^n}\\left[\\int_0^{\\tau_k}\\bar X_tdA^n_t \\right] = \\mathbb E^{\\mathbb Q^n}\\left[-(A^n_{-}\\cdot\\bar X)_{\\tau_k} +\\bar X_{\\tau_k}A^n_{\\tau_k}\\right] = \\mathbb E^{\\mathbb Q^n}\\left[\\bar X_{\\tau_k}A^n_{\\tau_k}\\right].\n\\end{displaymath}\nWhereas, in the other term in $T_2$, we get\n\\begin{displaymath}\n\\begin{array}{rcl}\n\\mathbb E^{\\mathbb Q^n}\\left[ \\int_0^{T}\\mathcal E^q_t dA^n_t\\right] &= &\\mathbb E^{\\mathbb Q^n}\\left[\\mathcal E^q_{T}A^n_{T} - (A^n_{-}\\cdot \\mathcal E^q)_{T}\\right] \\\\\n&= &\\mathbb E^{\\mathbb Q^n}\\left[\\mathcal E^q_{T} - (A^n_{-}\\cdot \\mathcal E^q)_{T}\\right] \\\\\n&= &\\mathbb E^{\\mathbb Q^n}\\left[((1- A^n_{-})\\cdot \\mathcal E^q)_{T}\\right] \\\\\n&= &\\mathbb E^{\\mathbb Q^n}\\left[(D^n\\cdot \\mathcal E^q)_{T}\\right].\\\\\n\\end{array}\n\\end{displaymath}\nusing integration by parts and localization, as in the proof of Lemma \\ref{6-16-1}, we can rewrite the latter expression as\n$$\n\\mathbb E\\left[\\int_0^T D^n_sZ^n_sq_se_sd\\kappa_s\\right].\n$$\nWe conclude that\n$$T_2 = -\\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[\\bar X_{\\tau_k}A^n_{\\tau_k}\\right] +\\mathbb E\\left[\\int_0^T D^n_sZ^n_sq_se_sd\\kappa_s\\right].$$\nCombining this with \\eqref{5214}, we obtain \n$$T_1 + T_2 = \\mathbb E\\left[\\int_0^T D^n_sZ^n_sq_se_sd\\kappa_s\\right] + \\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[A^n_{\\tau_k}\\left(x + \\int_0^{\\tau_k}\\widehat H_sdS_s +\\bar X_{\\tau_k} \\right)\\right] - \\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[A^n_{\\tau_k}\\bar X_{\\tau_k} \\right].$$\nAs both limits in the right-hand side exist and by positivity of $\\left(x + \\int_0^{\\tau_k}\\widehat H_sdS_s +\\bar X_{\\tau_k} \\right)$, established in \\eqref{5232}, we can bound the difference of the limits as\n\\begin{displaymath}\\begin{array}{c}\n \\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[A^n_{\\tau_k}\\left(x + \\int_0^{\\tau_k}\\widehat H_sdS_s +\\bar X_{\\tau_k} \\right)\\right] - \\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[A^n_{\\tau_k}\\bar X_{\\tau_k} \\right] \\\\\n \\leq \n \\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[x + \\int_0^{\\tau_k}\\widehat H_sdS_s +\\bar X_{\\tau_k}\\right] - \\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[A^n_{\\tau_k}\\bar X_{\\tau_k} \\right] \\\\\n \\leq \n \\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[x + \\int_0^{\\tau_k}\\widehat H_sdS_s\\right] + \n \\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[(1-A^n_{\\tau_k})\\bar X_{\\tau_k} \\right].\\\\\n\\end{array}\n\\end{displaymath}\nBy definition of $\\mathcal M'$, $\\bar X$ is a uniformly integrable martingale under $\\mathbb Q^n$. Therefore, as $(1-A^n_{\\tau_k})$ is bounded, we can pass the limit inside of the expectation to obtain $\\lim\\limits_{k\\to\\infty}\\mathbb E^{\\mathbb Q^n}\\left[(1-A^n_{\\tau_k})\\bar X_{\\tau_k} \\right] = 0$. In turn $x + \\int_0^{\\cdot}\\widehat H_sdS_s$ is a supermartingale under $\\mathbb Q^n$ (see the argument in the proof of Lemma \\ref{6-16-1}). We conclude that \n$$\nT_1 + T_2 \\leq x + \\mathbb E\\left[\\int_0^T D^n_sZ^n_sq_se_sd\\kappa_s\\right] = x + \\int_0^T q_s\\rho^n_sdK_s,\n$$\nfor some $\\rho^n$, which is well-defined by Corollary \\ref{corRzd}, and such that $y^n(1,\\rho^n)\\in z\\mathcal B(x,q)$, since $y^nZ^nD^n\\in z\\mathcal G(x,q)$. \nThus, from \\eqref{eqndefz} and \\eqref{5212}, we get\n\\begin{displaymath}\n\\begin{array}{rcl}\nz&\\leq &(x + \\int_0^T q_s\\rho^n_sdK_s)y^n -\\lim\\limits_{n\\to\\infty}y^n\\mathbb E\\left[Z^n_T \\int_0^T\\widehat V_tdA^n_t \\right].\\\\\n\\end{array}\n\\end{displaymath}\nBy optimality of $\\widehat Y$, $z\\geq (x + \\int_0^T q_s\\rho^n_sdK_s)y^n\\geq 0$. Therefore, by nonnegativity of $\\mathbb E\\left[Z^n_T \\int_0^T\\widehat V_tdA^n_t \\right]$, and since $y^n$ converges to a strictly positive limit, we conclude that\n$$\\lim\\limits_{n\\to\\infty}\\mathbb E\\left[Z^n_T \\int_0^T\\widehat V_tdA^n_t \\right] = 0.$$\nApplying integration by parts and localization we deduce the assertion of the lemma.\n\\end{proof}\n\\begin{Remark} We emphasize again that $dA^n$ are probability measures on $[\\theta _0,T]$, which can have mass at the endpoints, and \n $D^n_t=1-A^n_{t-}, \\theta _0\\leq t\\leq T$, $D^n_{T+}=0$.\n\\end{Remark}\n\nIn the subsequent part, we will follow the notations of Lemma \\ref{lem5211} and we will work with a further subsequence, still denoted by $n$, such that\n$y_n\\rightarrow y>0$, \n\\begin{equation}\n\\label{cone-product}\nZ^nD^n\\rightarrow \\frac{\\widehat Y}{y}, \\quad (d\\kappa\\times\\mathbb P)-a.e, \n\\end{equation} and\n\\begin{equation}\\label{fast}\n\\sum _{k=n}^{\\infty}\\mathbb E\\left[Z^n_T \\int_0^T\\widehat V_tdA^n_t \\right]\\leq \\frac{2^{-n}}{n},\n\\end{equation}\nwhere the existence of a subsequence satisfying \\eqref{fast} follows from \\eqref{86*}.\nThe next results is a Komlos-type lemma, largely based on the results in \\cite{CzichSchachOSSup}, applied to the \\emph{double-sequence} of processes $(Z^n, \\left (D^{n}\\right)^{-1})$. We observe that the process $\\left ( D^n\\right )^{-1}$ takes values in $[1, \\infty]$ is increasing, left-continuous, and satisfies\n$$\\left ( D^n _t\\right )^{-1}=1,\\quad 0\\leq t\\leq \\theta _0.$$\n\\begin{Lemma}\\label{5219} \nLet the conditions of Theorem \\ref{mainTheorem2} hold. In the notations of Lemma \\ref{lem5211},\nfor each $n$, there exist a finite index $N(n)$ and convex weights\n$$\\alpha _{n,k}\\geq 0,\\quad k=n,..., N(n),\\quad \\sum _{k=n}^{N(n)}\\alpha _{n,k}=1,$$ \nand there exists a strong optional super-martingale $\\widehat Z$ and a non-decreasing (not necessarily left-continuous) process $\\widehat D$ \nwith $\\widehat D_t=0, \\ \\forall 0\\leq t\\leq \\theta_0$, \nsuch that, simultaneously, \n\\begin{enumerate}\n\\item $\\tilde Z^n{:=} \\sum _{k=n}^{N(n)}\\alpha _{n,k}Z^k\\rightarrow \\widehat Z$ in the sense of \\cite{CzichSchachOSSup} i.e. for any stopping time $0\\leq \\tau\\leq T$ we have \n$$\\tilde Z^n _{\\tau} {\\longrightarrow}^{\\mathbb{P}} \\widehat Z_{\\tau},$$\nand \n\\item \n$$ \\mathbb{P}\\left (\\sum _{k=n}^{N(n)}\\alpha _{n,k}(D^k_t)^{-1}\\rightarrow (\\widehat D_t )^{-1}, \\quad \\forall 0\\leq t \\leq T+ \\right)=1$$\n\n We have set $D_{T+}=D^k_{T+}=0$. \\end{enumerate}\nFurthermore, with the notation\n\\begin{equation}\\label{a-hat}\n \\widehat A _t {=} 1-\\widehat D_{t+}, \\quad 0\\leq t\\leq T,\\quad \\widehat A_{0-}=0, \n\\end{equation}\nwe have the probability measure $d\\widehat A$ on $[\\theta _0,T]$ such that\n$$\\mathbb{P}\\left (\\widehat D_t\\in \\left[1-\\widehat A_t, 1-\\widehat A_{t-}\\right], \\ \\forall \\ \\theta _0 \\leq t\\leq T)\\right)=1.$$\nDenoting by \n\\begin{equation}\\label{a-tilde-n}\n\\tilde A_t^n{:=} 1- \\left (\\sum _{k=n}^{N(n)}\\alpha _{n,k}(1-A^k_t)^{-1}\\right)^{-1}=1-\n\\left (\\sum _{k=n}^{N(n)}\\alpha _{n,k}(D^k_{t+})^{-1}\\right)^{-1}, \\quad 0\\leq t\\leq T, \\quad \\tilde A ^n_{0-}=0,\n\\end{equation}\nthe point-wise convergence in time (at all times where there is continuity) additionally implies \n$$d\\tilde A^n\\rightharpoonup d\\widehat A,\\quad \\mathbb{P}-a.e.$$ in the sense of weak convergence of probability measures on $[\\theta _0, T]$. \\end{Lemma}\n\\begin{proof}\nThe proof reduces to applying the Komlos-type results in \\cite{CzichSchachOSSup} to the sequence $Z^n$ and \\cite[Proposition 3.4]{campi-schachermayer} to the sequence $(D^{n})^{-1}$, simultaneously. We observe that: \n\\begin{enumerate}\n\\item first, no bounds are needed for $D^{-1}$'s since we can apply the unbounded Komlos lemma in \\cite[Lemma A1.1]{DS} (and Remark 1 following it) to the proof from \\cite[Propositions 3.4]{campi-schachermayer}, and this works even for infinite values (according to Remark 1 after \\cite[Lemma A1.1]{DS}). Also, predictability can be replaced by optionality, without any change to the proof.\n\\item the Komlos arguments can be applied to both sequences $Z^n$ and $D^n$ simultaneously, with the same convex weights. In order to do this, we first apply Komlos to one sequence, then replace both original sequences by their convex combinations with the weights just obtained, and then apply Komlos again for the other sequence, and update the convex weight to both sequences.\n\\end{enumerate}\n\n\\end{proof}\n\n\\begin{Corollary} Let the conditions of Theorem \\ref{mainTheorem2} hold. In the notations of Lemma \\ref{5219}, we have the representation\n$$\\widehat Y= y\\widehat Z\\widehat D,\\quad (d\\kappa\\times\\mathbb P)-a.e.$$\n\\end{Corollary}\n\\begin{proof} Consider an observation that, for non-negative\nnumbers $a_k, b_k$, $k=n,\\dots, N(n)$, we have\n$$\\min _{k=n,...,N(n)} \\left (\\frac{a_k}{b_k}\\right)\\leq \\frac{\\sum _{k=n}^{N(n)}\\alpha _{n,k}a^k}{\\sum _{k=n}^{N(n)}\\alpha _{n,k}b_k}\\leq \\max _{k=n,...,N(n)} \\left (\\frac{a_k}{b_k}\\right).$$\\\nWe apply this to $a_k=Z^k, b_k=(D^k)^{-1},$ to obtain\n$$\\min _{k=n,...,N(n)} \\left (Z^kD^k\\right)\\leq \n \\frac{\\sum _{k=n}^{N(n)}\\alpha _{n,k}Z^k}{\\sum _{k=n}^{N(n)}\\alpha _{n,k}(D^k)^{-1}}=\\tilde Z^n \\tilde D^n \\leq\n \\max _{k=n,...,N(n)} \\left (Z^kD^k\\right),$$\npointwise a.e. in the product space, where $\\tilde D^n {:=} \\frac{1}{\\sum _{k=n}^{N(n)}\\alpha _{n,k}(D^k)^{-1}}$. \n Since $Z^nD^n\\rightarrow \\widehat Y\/ y$, $(d\\kappa\\times \\mathbb P)$-a.e., we conclude that both\n $\\left(\\min _{k=n,...,N(n)} \\left (Z^kD^k\\right)\\right)_{n\\in\\mathbb N}$ and $\\left(\\max _{k=n,...,N(n)} \\left (Z^kD^k\\right)\\right)_{n\\in\\mathbb N}$ converge to $\\frac{\\widehat Y}{y}$, $(d\\kappa\\times \\mathbb P)$-a.e.,\n therefore\n \\begin{equation}\n \\label{conv-product}\\tilde Z^n \\tilde D^n \\rightarrow \\frac{\\widehat Y}{y},\\quad (d\\kappa\\times \\mathbb P)-a.e.\n \\end{equation}\nUsing \\eqref{cone-product} above, if we did not plan to identify the limit $\\widehat Z$ as a strong-supermartingale, but only as a $(d\\kappa\\times \\mathbb P)$-a.e. limit in the product space, we could only apply Komlos arguments to the single sequence $D^{-1}$, to conclude convergence of the other convex combination $\\tilde Z^n$ to $\\frac{\\widehat Y}{y\\widehat D}$, defined up to a.e. equality in the product space.\n \n \n With our (stronger) double Komlos argument, we have that, in addition to \\eqref{cone-product} we have \\begin{equation}\n \\label{conv-tau}\\widehat Z_{\\tau}\\widehat D_{\\tau}=\\mathbb{P}-\\lim _n \\tilde Z^n_{\\tau}\\tilde D^n _{\\tau},\\quad \\textrm{for\\ every~[0,T]-valued~stopping \\ time \\ }\\tau,\n \\end{equation}\n (where $\\widehat Z$ is a strong supermartingale, and $\\widehat D$ is well defined at all times).\n\n It remains to prove that $\\widehat Y\/y=\\widehat Z\\widehat D$, $(d\\kappa\\times\\mathbb P)$-a.e. We point out that the convergence \\eqref{conv-tau} is topological. Let us consider an arbitrary optional set $O\\subset \\Omega \\times [0,T]$ and fix an upper bound $M$. It follows from \\eqref{conv-product} that\n $$1_O \\min \\left \\{\\tilde Z^n \\tilde D^n, M \\right \\}\\rightarrow 1_O \\min \\left \\{\\frac{\\widehat Y}{y}, M \\right \\},\\quad (d\\kappa\\times \\mathbb P)-a.e.\n$$\nTherefore, we get\n\\begin{equation}\\label{hatY}\\mathbb{E}\\left [\\int _0^T 1_O(t, \\cdot) \\min \\left \\{\\tilde Z^n_t \\tilde D^n _t, M \\right \\} d\\kappa_t\\right]\\rightarrow \n\\mathbb{E}\\left [\\int _0^T 1_O(t, \\cdot) \\min \\left \\{\\frac{\\widehat Y_t}{y}, M \\right \\} d\\kappa_t\\right].\n\\end{equation} \nRecall that the stochastic clock has a density $d\\kappa_t=\\varphi _t dK_t$ with respect to the deterministic clock $dK _t$. For each fixed $t$, from \\eqref{conv-tau} we have\n$$1_O(t, \\cdot) \\min \\left \\{\\tilde Z^n_t \\tilde D^n _t, M \\right \\} \\varphi _t\\rightarrow 1_O(t, \\cdot) \\min \\left \\{\\widehat Z_t \\widehat D _t, M \\right \\} \\varphi _t,\\quad \\textrm{in}-\\mathbb{P}.$$\nRecalling that $\\mathbb{E}[\\varphi _t]<\\infty$ we have, for fixed $t$, that \n$$ M \\times \\mathbb{E}[\\varphi _t]\\geq \\mathbb{E}\\left [1_O(t, \\cdot) \\min \\left \\{\\tilde Z^n_t \\tilde D^n _t, M \\right \\} \\varphi _t \\right ]\\rightarrow \\mathbb{E}\\left [1_O(t, \\cdot) \\min \\left \\{\\widehat Z_t \\widehat D _t, M \\right \\} \\varphi _t\\right ].\n$$\nWe integrate the above with respect to the deterministic clock $dK_t$ to obtain\n\\begin{equation}\n\\begin{split}\n\\mathbb{E}\\left [\\int _0^T 1_O(t, \\cdot) \\min \\left \\{\\tilde Z^n_t \\tilde D^n _t, M \\right \\} d\\kappa_t\\right] =\\int _0^T \\mathbb{E}\\left [1_O(t, \\cdot) \\min \\left \\{\\tilde Z^n_t \\tilde D^n _t, M \\right \\} \\varphi _t \\right ]dK_t\\rightarrow \\\\\n\\rightarrow \\int _0^T \\mathbb{E}\\left [1_O(t, \\cdot) \\min \\left \\{\\widehat Z_t \\widehat D _t, M \\right \\} \\varphi _t \\right ]dK_t=\n\\mathbb{E}\\left [\\int _0^T 1_O(t, \\cdot) \\min \\left \\{\\widehat Z_t \\widehat D _t, M \\right \\} d\\kappa_t\\right]\\end{split}\n\\end{equation}\nTogether with \\eqref{hatY} we have\n$$\\mathbb{E}\\left [\\int _0^T 1_O(t, \\cdot) \\min \\left \\{\\frac{\\widehat Y_t}{\\hat y}, M \\right \\} d\\kappa_t\\right]=\\mathbb{E}\\left [\\int _0^T 1_O(t, \\cdot) \\min \\left \\{\\widehat Z_t \\widehat D _t, M \\right \\} d\\kappa_t\\right],$$\nwhich holds for any optional set $O$ and any bound $M$, therefore $\\widehat Y= y\\widehat Z\\widehat D$, $(d\\kappa\\times \\mathbb P)$-a.e.\n \\end{proof}\n\\begin{proof}[Proof of the Theorem \\ref{ExactSlackness}]\nSince$$\\sum _{k=n}^{N(n)}\\alpha _{n,k}(D^k_t)^{-1}\\rightarrow (\\widehat D_t )^{-1}, \\quad t\\in[0,T], $$\nrecalling that $\\widehat A$ was defined from $\\widehat D$ and \\eqref{a-tilde-n}. \nAs the processes $A^n$, therefore $\\tilde A^n$ only increase by jumps.\nTherefore\nwe get\n \\begin{equation}\\label{11291}\n \\begin{array}{rcl}\n \\Delta \\tilde A^n_s &= & \\frac{1}{\\sum\\limits_{k=n}^{N(n)}\\alpha _{n,k}\\frac{1}{D^k_{s}}} - \\frac{1}{\\sum\\limits_{k=n}^{N(n)}\\alpha _{n,k}\\frac{1}{D^k_{s+}}}\\\\\n &= & \\sum\\limits_{k=n}^{N(n)}\\underbrace{\\frac{\\alpha _{n,k}\\frac{1}{D^k_{s+}}}{\\sum\\limits_{k=n}^{N(n)}\\alpha _{n,k}\\frac{1}{D^k_{s+}}}}_{\\leq 1}\\frac{\\Delta A^k_s}{D^k_{s}} \n \\underbrace{\\frac{1}{\\sum\\limits_{k=n}^{N(n)}\\alpha _{n,k}\\frac{1}{D^k_{s}}}}_{\\leq 1}.\\\\\n \\end{array}\n \\end{equation}\n As $\\frac{\\alpha _{n,k}\\frac{1}{D^k_{s+}}}{\\left(\\sum\\limits_{k=n}^{N(n)}\\alpha _{n,k}\\frac{1}{D^k_{s+}}\\right)}\\leq 1$ and $\\sum\\limits_{k=n}^{N(n)}\\alpha _{n,k}\\frac{1}{D^k_{s}} \\geq 1$, we can bound the latter term in \\eqref{11291} by $\\sum\\limits_{k=n}^{N(n)}\\frac{\\Delta A^k_s}{D^k_{s}}$ for $s\\in[0,T]$.\n We deduce that\n $$\\Delta \\tilde A^n_s \\leq \\sum\\limits_{k=n}^{N(n)}\\frac{\\Delta A^k_s}{D^k_{s}},\\quad s\\in[0,T].$$ \nTherefore, we have\n$$\\int_{0}^{t} \\widehat V_u d\\tilde A^n _u \\leq \\sum _{k=n}^{N(n)}(D^k_{t})^{-1} \\int_{0}^{t} \\widehat V_udA^k_u,\\quad t\\in[0,T],$$\nand thus, we obtain\n\\begin{equation}\\label{main-inez}\n\\left(\\min _{k=n, \\dots, N(n)} (Z^k_tD^k_t) \\right) \\int_{0}^{t} \\widehat V_u d\\tilde A^n _u \\leq \\sum _{k=n}^{N(n)}(Z^k_tD^k_t)(D^k_{t})^{-1} \\int_{0}^{t} \\widehat V_udA^k_u= \\sum _{k=n}^{N(n)}Z^k_t\\int_{0}^{t} \\widehat V_udA^k_u,\\quad t\\in[0,T].\n\\end{equation}\nSince the process\n$$L^n_t:=\\sum _{k=n}^{N(n)}Z^k_t\\int_{0}^{t} \\widehat V_udA^k_u, \\quad 0\\leq t\\leq T$$\nis a non-negative right-continuous submartingale, the maximal inequality and \\eqref{fast} together imply\n $$\\mathbb{P}\\left (\\sup _{0\\leq t\\leq T}L^n_t\\geq \\frac 1n \\right )\\leq n\\mathbb{E}[L^n_T]\\leq 2^{-n},\n$$\nso \n\\begin{equation}\n\\label{l}\\sup _{0\\leq t\\leq T}L^n_t\\rightarrow 0,\\quad \\mathbb{P}-a.s.\n\\end{equation}\nSince $Z^nD^n \\rightarrow \\frac{1}{y}\\widehat Y>0$, $(d\\kappa\\times \\mathbb P)$-a.e., consequenlty \n$$\\min _{k=n, \\dots, N(n)} (Z^kD^k)\\rightarrow \\frac{1}{y}\\widehat Y>0\\quad \\left(d\\kappa \\times \\mathbb{P}\\right)-a.e.,$$\nwe obtain from\\eqref{main-inez} and \\eqref{l} that \nthe increasing RC process\n$$\\tilde L^n_t:=\\int _0^t \\widehat V_u d\\tilde A^n_u,\\quad 0\\leq t\\leq T,$$\nconverges to zero in the product space. Denoting by $\\mathcal{O}\\subset \\Omega\\times [0,T]$ the exceptional set where convergence to zero does not take place, and taking into account that $\\tilde L^n$ are increasing, we have that for \nfor $$T^{\\mathcal{O}}(\\omega)=\\inf \\{0\\leq t\\leq T: (\\omega ,t)\\in \\mathcal{O}\\},$$\nwe have \n$$(T^{\\mathcal{O}},T]\\subset \\mathcal{O}.$$\nNow \n$$(d\\kappa\\times \\mathbb P)(\\mathcal{O})=0, $$\nimplies that \n$T^{\\mathcal{O}}\\geq T,~\\mathbb{P}-a.s.$ (here used an assumption that $T$ is the minimal time horizon in the sense that the deterministic clock $K$ is such that $K_t R_1$, concentric with $D_1$, is the image of the annulus $D_2\\backslash D_1$ by a radial transformation. In cylindrical coordinates, this transformation is given by:\n\n\\begin{equation}\\begin{cases} r = (r'-R_1) R_2\/(R_2-R_1) \\; \\mathrm{for} \\; R_1 \\leq r'\\leq R_2, \\; \\\\ \\theta = \\theta', \\; z = z'. \\;\n\\label{pendrysmap}\\end{cases}\\end{equation}\n\nAs for the outside of the disk $D_2$, the map between the two copies of $\\mathbb{R}^2\\backslash D_2$ is the identity map.\n\nThe material properties given by rule (\\ref{equivalence_rule}) corresponding to this transformation provide an ideal invisibility cloak: outside $D_2$, everything behaves as if we were in free space, including the propagation of electromagnetic waves across the cloak, and is completely independent of the content of $D_1$.\n\n\nNow, rule (\\ref{equivalence_rule}) may be applied to $D_2$ containing objects with arbitrary electromagnetic properties so that a region cloaked by this device is still completely hidden but has the appearance of the objects originally in $D_2$. We may call this optical effect masking \\cite{teixeira} or ``polyjuice'' effect.\n\n\n\n\n\n\n\n\n\n\n\n\n\\begin{figure}[h\n{\\centering \\includegraphics[width=0.49\\textwidth]{juicemesh2.pdf}\n\\caption{\\label{mesh} This figure shows a part of the triangular mesh used for the finite element modeling of the scattering problem of Fig. \\ref{fig2}. The singular behavior of the permittivity and of the permeability requires a very fine mesh along the inner boundary of the cloak in order to achieve a satisfactory accuracy with the numerical model. }}\n\\end{figure}\n\n\n\\begin{figure}[h\n{\\centering \\includegraphics[width=0.49\\textwidth]{etopim2.jpg}\n\\caption{\\label{fig1} A conducting triangular cylinder is scattering cylindrical waves.}}\n\\end{figure}\n\n\n\\begin{figure}[h\n{\\centering \\includegraphics[width=0.49\\textwidth]{etopim1.jpg}\n\\caption{\\label{fig2} A triangular cylinder different from the one on Fig. \\ref{fig1} is surrounded by a cloak designed to reproduce the scattering pattern of the Fig. \\ref{fig1} triangular cylinder in spite of the change of diffracting object. Of course, the diffracting object inside the cloak may be arbitrary as far as it is small enough to fit inside the cloak}}\n\\end{figure}\n\n\n\\begin{figure}[h\n{\\centering \\includegraphics[width=0.49\\textwidth]{compare.pdf}\n\\caption{\\label{fig3} The value of the electric field (the real part of $E_z$) on a circle of radius $4\\lambda$ concentric with the cloak is represented as a function of the position angle $\\theta$ (increasing counterclockwise and with $\\theta=0$ corresponding to the point the most on the right). The three configurations considered here are the ones of Fig. \\ref{fig2} (coated), Fig. \\ref{fig1} (original), and the triangle of Fig. \\ref{fig2} without the coating (reversed). } }\n\\end{figure}\n\n\nFigs. \\ref{fig1} and \\ref{fig2} show the effect of masking on a scattering structure. On Fig. \\ref{fig1}, a cylindrical TM wave emitted by a circular cylindrical antenna is scattered by a conducting triangular cylinder (the longest side of the cross section is $1.62\\lambda$ and $\\varepsilon_r=1 + 40 i $). The field map represents the longitudinal electric field $E_z(x,y)$ and the outer boundary of the cloak is shown to ease the comparison with the masked case. On Fig. \\ref{fig2}, the same cylindrical TM wave is scattered by a masked triangular cylinder (but the diffracting object inside the cloak may be arbitrary as far as it is small enough to fit inside the cloak). This triangular cylinder is the symmetric of the previous one with respect to the horizontal plane containing the central fibre of the cylindrical antenna. This bare scatterer would therefore give the Fig. \\ref{fig1} image inverted upside-down but, here, this object is surrounded by a cloak in order to give the very same scattering as before. Indeed, on both sides, the electric fields outside the cloak limit are alike.\n\n Fig. \\ref{fig3} highlights the different scattering patterns by displaying the value of $\\Re e(E_z)$ on a circle of radius $4 \\lambda$ located around the antenna-scatterer system in the three following cases: the case of Fig. \\ref{fig1} (original) with the triangle alone, the case of Fig. \\ref{fig2} (coated), and the triangle of Fig. \\ref{fig2} without the coating (reversed). It is obvious that the coating restores the field distribution independently of the object present in the central hole.\n\n\nThe numerical computation is performed using the finite element method (via the free GetDP \\cite{getdp} and Gmsh \\cite{gmsh} software tools). The mesh is made of 148,000 second order triangles including the Perfectly Matched Layers used to truncate the computation domain . The singularity of $\\varepsilon$ and $\\mu$ requires a very fine mesh in the vicinity of the inner boundary of the cloak (see Fig. \\ref{mesh}) and is also responsible for the small discrepancies between the numerical model and a perfect cloak (see Fig. \\ref{fig2}) --- including the non zero field in the hole of the cloak.\n\nNote that a small technical problem arises in practice when rule (\\ref{equivalence_rule}) is applied: the material properties are defined piecewise on various domains and it is very useful to know explicitly the boundaries of these domains, e.g. to build the finite element mesh (see Fig. \\ref{mesh}). These boundaries are curves in the cross section and are thus contravariant objects. Therefore, their transformation requires the inverse map $\\varphi^{-1}$ from $M$ to $N$. Fortunately, map (\\ref{pendrysmap}) is very simple to invert.\n\nMore explicitly, for a given curve $\\mathbf{x}(t)$ of parameter $t$ in the initial Cartesian coordinates, its push forward by Pendry's map is:\n\\begin{equation}\n\\mathbf{x}'(t)=\\varphi^{-1}(\\mathbf{x}(t))=(\\frac{R_2-R_1}{R_2}+ \\frac{R_1}{\\|\\mathbf{x}(t) \\|})\\mathbf{x}(t),\n\\end{equation}\nwith the same variation of the parameter $t$. Note that the most common curves used in the design of devices, i.e. line segments and arc of circles, are transformed to less usual curves except for radial segments (with respect to the center of the cloak) and arc of circles concentric with the cloak.\n\n On Fig. \\ref{fig2}, the image by $\\varphi^{-1}$ of the triangle of Fig. \\ref{fig1} is the curvilinear triangle inside the coating region of the cloak. In practice, this anamorphosis of the triangle is described by three splines interpolating each 40 points that are images of points of the segments by $\\varphi^{-1}$.\n\n\n\n\nTransformation optics do not offer only the possibility make optically disappear objects in invisibility cloaks but also to completely tune their optical signature i.e. to give them an arbitrary appearance.\n The possibility to place an object inside the coating of the cloak and that it will therefore appear different was already considered in \\cite{opl} where a point source was shifted creating a mirage effect. A more general case is considered here since both the shape and the position of the object placed in the coating are modified .\n Note that if the object used to create the illusion is perfectly conducting and surrounds the central point of the cloak, its anamorphosis will surround the inner boundary of the cloak and it \\emph{de facto} suppresses the singular behavior of the material properties. In this particular case, the present approach is similar to the the hiding under the carpet idea of Li and Pendry \\cite{carpet} but here a bounded object replaces the infinite reflecting plane. This technique can be naturally extended to cloaks of arbitrary shapes \\cite{arbitrary}.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nLet $\\I\\subset\\R$ be an interval and $f\\colon\\I\\to\\R$. The \\emph{difference operator} is given by\n\\[\n \\Delta_hf(x)=f(x+h)-f(x)\n\\]\nfor any $x\\in\\I$ and $h\\in\\R$ \\st\\ $x+h\\in\\I$. In 1954 Wright~\\cite{Wri54} considered the class of functions~$f:\\I\\to\\R$, for which the function $\\Delta_h f(\\cdot)$ is non-decreasing for any $h>0$ small enough to guarantee that all the involved arguments belong to~$\\I$. This monotonicity property is equivalent to the condition\n\\begin{equation}\\label{eq:Wright}\n f\\bigl(tx+(1-t)y\\bigr)+f\\bigl((1-t)x+ty\\bigr)\\xle f(x)+f(y)\\quad \\big(x,y\\in\\I,\\;t\\in[0,1]\\big).\n\\end{equation}\nNowadays a~function fulfilling~\\eqref{eq:Wright} is called to be \\emph{Wright-convex}. For $t=\\frac{1}{2}$ we get immediately\n\\[\n f\\Bigl(\\frac{x+y}{2}\\Bigr)\\xle\\frac{f(x)+f(y)}{2}\\quad\\bigl(x,y\\in\\I\\bigr),\n\\]\nso any Wright-convex function is necessarily Jensen-convex. It was shown by Ng~\\cite{Ng87} that any Wright-convex function is a~sum of the convex function (in the usual sense) and the additive one. Such representation leads to the simple (and now classical) proof that the inclusion in question is proper. Namely, if $a\\colon\\R\\to\\R$ is a~discontinuous additive function and $f(x)=|a(x)|$, then $f$~is discontinuous and Jensen-convex. Hence the graph of~$f$ is not dense on the whole plane. But it is well-known (\\cf~\\eg~\\cite{Kuc09}) that the graphs of discontinuous additive functions (and also -- by Ng's representation -- the graphs of discontinuous Wright-convex functions) mapping~$\\R$ into~$\\R$ are dense on the whole plane. That is why~$f$ is not Wright-convex.\n\\par\\medskip\nRecently Nikodem, Rajba and W\u0105sowicz in the paper~\\cite{NikRajWas12JMAA} considered the related comparison problem for the classes of $n$-Wright-convex functions and $n$-Jensen-convex functions (the explanation is given below in this section). If $n\\in\\N$ is odd, they have shown that the first class is the proper subclass of the latter one. The example of $n$-Jensen-convex function which is not $n$-Wright-convex was $f\\colon\\R\\to\\R$ \\st\\ $f(x)=\\bigl(\\max\\{a(x),0\\}\\bigr)^n$, where $a\\colon\\R\\to\\R$ is the certain discontinuous additive function. Next P\\'ales~\\cite{Pal13} proved that this example works for any discontinuous additive function. The authors of~\\cite{NikRajWas12JMAA} did not solve the problem, if~$n$ is even. The objective of the present paper is to give the complete solution by showing that the inclusion is proper for all $n\\in\\N$. Moreover, the construction of our example does not depend on the parity of~$n$. Then for odd~$n$ we have a~solution different from the ones given in the papers~\\cite{NikRajWas12JMAA,Pal13}.\n\\par\\medskip\nLet us recall the concepts of higher-order Jensen-convex (coming from Popoviciu~\\cite{Pop34}) and Wright-convex functions (introduced by Gil\\'anyi and P\\'ales~\\cite{GilPal01}).\n\\par\\medskip\nThe iterates of ~difference operator are we defined as usual: for $n\\in\\N$\n\\[\n \\Delta_{h_1\\dots h_nh_{n+1}}f(x)=\\Delta_{h_1\\dots h_n}\\bigl(\\Delta_{h_{n+1}}f(x)\\bigr)\n\\]\nprovided that all the involved arguments belong to~$\\I$. If $h_1=\\dots=h_n=h$, then we write \n\\[\n \\Delta_h^nf(x)=\\Delta_{h\\dots h}f(x).\n\\]\nIn particular, $\\Delta_h^1f(x)=\\Delta_hf(x)$. \nLet $n\\in\\N$ and $f\\colon\\I\\to\\R$ be a~function. We say that $f$ is \\emph{Wright-convex of order~$n$} (or $n$-\\emph{Wright-convex} for short), if\n\\[\n \\Delta_{h_1\\dots h_{n+1}}f(x)\\xge 0\n\\]\nfor any $x\\in\\I$ and $h_1,\\ldots,h_{n+1}>0$ \\st\\ $x+h_1+\\dots+h_{n+1}\\in\\I$. If the above inequality is required only for $h_1=\\dots=h_{n+1}=h>0$, precisely, if\n\\[\n \\Delta_h^{n+1}f(x)\\xge 0\n\\]\nfor all $x\\in I$ and $h>0$ such that $x+(n+1)h\\in\\I$, then~$f$ is called \\emph{Jensen-convex of order~$n$} ($n$-\\emph{Jensen-convex} for short). It is easy to see that $1$-Wright-convexity and $1$-Jensen-convexity are equivalent to Wright-convexity and Jensen-convexity, respectively.\n\\par\\medskip\nBy the above definitions any $n$-Wright-convex function is necessarily $n$-Jensen-convex. As we announced above, in this paper we show that this inclusion is proper.\n\n\\section{Main result}\nLet us start with two preparatory notes. The following equations hold for a~power function function $f(x)=x^n$ (\\cf~\\cite[Lemma 15.9.2]{Kuc09}):\n\\begin{equation}\\label{eq:delta-xn}\n \\Delta_{h_1\\ldots h_n}\\bigl(x^n\\bigr)=n!\\cdot\\prod_{i=1}^nh_i \\qquad\\text{and}\\qquad \\Delta_h^n\\bigl(x^n\\big)=n!\\cdot h^n.\n\\end{equation}\n\\par\\medskip\nThe formula below was given in~\\cite{NikRajWas12JMAA} (as a~part of the proof of Corollary~2.2).\n\\begin{lem}\\label{lemat2}\nLet $a\\colon\\R\\to\\R$ be an additive function, $g\\colon\\R\\to\\R$ and $n\\in\\N$.\nThen\n\\[\n \\Delta_{h_1\\dots h_n}(g\\circ a) (x)=\\Delta_{a(h_1)\\dots\\,a(h_n)}g\\bigl(a(x)\\bigr)\n\\]\nfor any $x,h_1,\\dots,h_n\\in\\R$.\n\\end{lem}\nNow we are in a position to state our main result.\n\\begin{thm}\\label{th_main}\nFor any $n\\in\\N$ the class of $n$-Wright-convex functions is a~proper subclass of the class of $n$-Jensen-convex functions.\n\\end{thm}\n\n\\begin{proof}\n Because every $n$-Wright-convex function is $n$-Jensen-convex, it is enough to construct an~$n$-Jensen-convex function $f\\colon\\R\\to\\R$ which is not $n$-Wright convex. To this end consider a~Hamel basis $H$ of the linear space $\\R$ over~$\\Q$ \\st~$h_0=1\\in H$. Any $x\\in\\R$ is uniquely represented as a linear combination\n\\[\n x=\\lambda_0h_0+\\sum_{t}\\lambda_th_t,\n\\]\nwhere $\\lambda_0,\\lambda_t\\in\\Q$, $h_t\\in H$. Obviously the functions $\\alpha(x)=\\lambda_0h_0$ and $\\beta(x)=x-\\alpha(x)$ are additive. Define the function $f\\colon\\R\\to\\R$ by\n\\[\n f(x)=\\bigl(\\alpha(x)\\bigr)^{n+1}+\\bigl(\\beta(x)\\bigr)^{n+1}.\n\\]\nSince the difference operator is additive, we infer by Lemma~\\ref{lemat2} and the equation~\\eqref{eq:delta-xn} that\n\\begin{multline}\\label{eq:dj}\n\\Delta_{h_1\\dots h_{n+1}}f(x)=\\Delta_{\\alpha(h_1)\\ldots\\alpha(h_{n+1})}\\bigl(\\alpha(x)\\bigr)^{n+1}+\\Delta_{\\beta(h_1)\\ldots\\beta(h_{n+1})}\\bigl(\\beta(x)\\bigr)^{n+1}\\\\\n=(n+1)!\\biggl(\\prod_{i=1}^{n+1}\\alpha(h_i)+\\prod_{i=1}^{n+1}\\beta(h_i)\\biggr).\n\\end{multline}\nIn particular,\n\\[\n \\Delta_h^{n+1}f(x)=(n+1)!\\Bigl(\\bigl(\\alpha(h)\\bigr)^{n+1}+\\bigl(\\beta(h)\\bigr)^{n+1}\\Bigr).\n\\]\nLet $h>0$. Because of the representation $h=\\alpha(h)+\\beta(h)$, it is easy to see that $\\bigl(\\alpha(h)\\bigr)^{n+1}+\\bigl(\\beta(h)\\bigr)^{n+1}>0$. Therefore $\\Delta_h^{n+1}f(x)>0$ which shows that~$f$ is $n$-Jensen-convex.\n\nTo prove that $f$ is not $n$-Wright-convex take now $h_1=-1+\\sqrt{2}$, $h_2=1$ and (if $n\\ge 2$) $h_3=\\dots=h_{n+1}=1+\\sqrt{2}$. Then $h_i>0$, $i=1,\\dots,n+1$. Moreover $\\alpha(h_1)=-1$, $\\alpha(h_i)=1$, $i=2,\\dots,n+1$ and $\\beta(h_2)=0$. By~\\eqref{eq:dj} we obtain \n\\[\n \\Delta_{h_1\\ldots h_{n+1}}f(x)=-(n+1)!<0.\n\\]\nIt means that~$f$ is not $n$-Wright-convex.\n\\end{proof}\n\n\\section{The classes of higher order strongly Jensen-convex functions and Wright-convex functions}\n\nLet $n\\in\\N$ and $c>0$. A~function $f\\colon\\I\\to\\R$ is called \\emph{strongly Wright-convex of order $n$ with modulus $c$} (or strongly $n$-Wright-convex with modulus $c$) if\n\\[\n\\Delta_{h_1\\ldots h_{n+1}}f(x)\\xge c (n+1)!\\,h_1\\ldots h_{n+1}\n\\]\nfor all $x\\in I$ and $h_1,\\ldots,h_{n+1}>0$ such that $x+h_1+\\dots+h_{n+1}\\in\\I$. If the above inequality is required only for $h_1=\\dots=h_{n+1}=h>0$, precisely, if\n\\[\n \\Delta_{h}^{n+1}f(x)\\xge c(n+1)!\\,h^{n+1}\n\\]\nfor all $x\\in I$ and $h>0$ such that $x+(n+1)h\\in\\I$, then $f$ is called \\emph{strongly Jensen-convex of order $n$ with modulus~$c$} (or strongly $n$-Jensen-convex with modulus $c$) (\\cf~ \\cite{GerNik11}).\nFor $c=0$ we arrive at standard $n$-Wright-convex and $n$-Jensen-convex functions, respectively.\n\\par\\medskip\nThe following characterization of strongly $n$-Wright-convex functions was recently given in~\\cite{GilMerNikPal15}.\n\\begin{thm}\\label{tw4}\nLet $n\\in\\N$ and $c>0$. A function $f\\colon\\I\\to\\R$ is strongly $n$-Wright-convex with modulus~$c$ if and only if the function $g\\colon\\I\\to\\R$ \\st\\ $g(x)=f(x)-cx^{n+1}$, is $n$-Wright-convex.\n\\end{thm}\nIt is surprising that so far nobody wrote explicitly the similar characterization of strongly $n$-Jensen-convex functions. We fill this gap now.\n\\begin{thm}\\label{tw5}\nLet $n\\in\\N$ and $c>0$. A function $f\\colon\\I\\to\\R$ is strongly $n$-Jensen-convex with modulus~$c$ if and only if the function $g\\colon\\I\\to\\R$ \\st\\ $g(x)=f(x)-cx^{n+1}$, is $n$-Jensen-convex.\n\\end{thm}\n\\begin{proof}\nSuppose first, that the function $f$ is strongly $n$-Jensen-convex with modulus $c$. Applying~\\eqref{eq:delta-xn} to $g$ we infer that\n\\[\n\\Delta_{h}^{n+1}g(x)= \\Delta_{h}^{n+1}f(x) - \\Delta_{h}^{n+1}(c x^{n+1}) \\xge c (n+1)!\\,h^{n+1}- c (n+1)!\\,h^{n+1}=0,\n\\]\nwhence $g$ is $n$-Jensen-convex.\n\nConversely, if $g$ is an $n$-Jensen-convex with modulus~$c$, use~\\eqref{eq:delta-xn} for $f(x)=g(x)+cx^{n+1}$:\n\\[\n\\Delta_{h}^{n+1}f(x)= \\Delta_{h}^{n+1}g(x) + \\Delta_{h}^{n+1}(c x^{n+1}) \\xge 0+ c (n+1)!\\,h^{n+1}=c (n+1)!\\,h^{n+1}.\n\\]\nThis proves that~$f$ is strongly $n$-Jensen-convex.\n\\end{proof}\nThe following result we derive from Theorems \\ref{th_main}, \\ref{tw4}, \\ref{tw5}.\n\\begin{thm}\nLet $n\\in\\N$ and $c>0$. The class of strongly $n$-Wright-convex functions with modulus~$c$ is a~proper subclass of the class of strongly $n$-Jensen-convex functions with modulus~$c$.\n\\end{thm}\n\\begin{proof}\nIt is enough to show, that there exists a strongly $n$-Jensen--convex function with modulus $c$, which is not strongly $n$-Wright--convex with modulus $c$. By Theorem~\\ref{th_main} there exists a function~$g$, which is $n$-Jensen-convex and $g$ is not $n$-Wright-convex. Let $f(x)=g(x)+cx^{n+1}$. Since $g$ is $n$-Jensen-convex, it follows by Theorem~\\ref{tw5} that $f$~is strongly $n$-Jensen-convex with modulus~$c$. By Theorem~\\ref{tw4}~$f$ is not strongly $n$-Wright-convex with modulus~$c$ (because the function $g(x)=f(x)-cx^{n+1}$ is not $n$-Wright-convex). This concludes the proof.\n\\end{proof}\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\n\n\n\n\\section{Introduction}\n\\label{se:intro}\n\nAfter the discovery of a Higgs boson at the Large Hadron Collider (LHC) \\cite{Aad:2012tfa,Chatrchyan:2012ufa} at CERN, the complete identification of this particle is ongoing. The properties of the discovered particle, such as its couplings, are determined experimentally in order to fully identify its nature. For the endeavour of the identification of this particle, \ninput from the theory side is needed in form of precise predictions for the production and decay processes in the Standard Model (SM) as well as in its extensions that are to be tested. It is also crucial to provide reliable uncertainty estimates of the theoretical predictions. \nUnderestimating this uncertainty might lead to wrong conclusions.\nIn the SM, predictions and error estimates are well advanced, and in SM extensions they are\nconsolidating as well (see, e.g., the reviews in \nRefs.~\\cite{Dittmaier2011,Dittmaier2012,Dittmaier:2012nh,Heinemeyer:2013tqa,deFlorian:2016spz,%\nSpira:2016ztx}).\n\nOne of the simplest extensions of the SM \nis the Two-Higgs-Doublet Model (THDM) \\cite{Lee:1973iz,HHGuide:1990} where a second Higgs doublet is added \nto the SM field content. The underlying gauge group \n$\\mathrm{SU}(3)_{\\text{C}} \\times \\mathrm{SU}(2)_{\\text{W}} \\times \\mathrm{U}(1)_{\\text{Y}}$ as well as the fermion content of the SM are kept. After spontaneous symmetry breaking, there are five physical Higgs bosons where three of them are neutral and two are charged. In the CP-conserving case, which we consider, one of the neutral Higgs bosons is CP-odd and two are CP-even with one of them being SM-like. \n\nEven such a simple extension of the SM \ncan help solving some questions that are unanswered in the SM. For example, CP-violation in the Higgs sector could provide solutions to the problem of baryogenesis~\\cite{Turok:1990zg,Davies:1994id,Cline:1995dg,Fromme:2006cm}, and inert THDMs contain a dark matter candidate \\cite{PhysRevD.18.2574,Dolle:2009fn}. An even larger motivation comes from the embedding of the THDM into more complex models, such as axion \\cite{Kim:1986ax, Peccei:1977hh} or supersymmetric models \\cite{Haber:1984rc}. Some of the latter are promising candidates for a fundamental theory, and supersymmetric Higgs sectors contain a THDM (in which the doublets have opposite hypercharges). Even though the THDM is unlikely to be the fundamental theory of nature, it provides a rich phenomenology, which can be used in the search for a non-minimal Higgs sector without being limited by constraints from a \nmore fundamental theory.\n\nIn this sense, it is obvious that the THDM should be tested against data,\nand phenomenological studies have been performed recently, e.g., \n{in Refs.~\\cite{Celis:2013rcs,Celis:2013ixa,Chang:2013ona,Harlander:2013mla,%\nHespel:2014sla,Baglio:2014nea,Broggio:2014mna,%\nBernon:2015qea,Bernon:2015wef,%\nGoncalves:2016qhh,Dorsch:2016tab,%\nCacciapaglia:2016tlr,Cacchio:2016qyh,Aggleton:2016tdd,Han:2017pfo}.} \nIn order to provide precise predictions within this model, not only leading-order \n(LO), but also next-to-leading-order (NLO) contributions have to be taken into account. For the calculation of NLO contributions, a proper definition of a renormalization scheme is mandatory. There is no unique choice, and applying different renormalization schemes can help to estimate the theoretical uncertainty of the prediction that originates from the truncation of the perturbative series. \nThe renormalization of the THDM has already been tackled in several publications:\nFirst, in Ref.~\\cite{Santos:1996vt}, the fields and masses were renormalized in the on-shell scheme, \nhowever, the prescription given there does not cover all parameters.\nIn Refs.~\\cite{Kanemura:2004mg,LopezVal:2009qy}, a minimal field renormalization was applied, and the field mixing conditions were used to fix some of the mixing angles. \nIn view of an automation of NLO predictions within the THDM a tool was written by \nDegrande~\\cite{Degrande:2014vpa} \nwhere all finite rational terms and all\ndivergent terms are computed using on-shell conditions or conditions within the ``modified minimal subtraction scheme'' ($\\overline{\\mathrm{MS}}$). \nThough automation is very helpful, often specific problems occur depending on the model, the process, or the renormalization scheme considered, and, it might be necessary to solve these ``manually''. Specifically, spontaneously broken gauge theories with extended scalar sectors pose issues with the renormalization of vacuum expectation values and the related ``tadpoles'', \njeopardizing gauge independence and perturbative stability in predictions.\nRenormalization schemes employing a gauge-independent treatment of the tadpole terms were described recently in Refs.~\\cite{Krause:2016oke,Denner:2016etu,Krause:2016xku}. \n\nIn this paper, we perform the renormalization of various types of THDMs \n(Type I, Type II, ``lepton-specific'', and ``flipped''),\ndescribe four different renormalization schemes\n(for each type), \nand provide explicit results facilitating their application in NLO calculations. The comparison of results obtained in these renormalization schemes allows for checking their perturbative consistency, i.e.\\ whether the expansion point for the perturbation series is chosen well and no unphysically large corrections are introduced. Knowing in which parts of the parameter space a renormalization scheme leads to a stable perturbative behaviour is important for the applicability of the scheme. \nIn addition, we investigate the dependence on the renormalization scale $\\mu_\\mathrm{r}$ which is introduced by defining some parameters via {$\\overline{\\mathrm{MS}}$}{} conditions. In order to investigate the $\\mu_\\mathrm{r}$ dependence consistently, we solve the renormalization group equations (RGEs) and include the running effects. We also make contact to different renormalization schemes suggested in the literature, including the recent formulations \\cite{Krause:2016oke,Denner:2016etu} \nwith gauge-independent treatments of tadpoles.%\n\\footnote{In our work we do not consider the ``tadpole-pinched'' scheme suggested\nin \\citere{Krause:2016oke}.\nFollowing the arguments of \\citeres{Denner:1994nn,Denner:1994xt}\nwe consider the ``pinch technique'' just as one of many physically\nequivalent choices to fix the gauge arbitrariness in off-shell\nquantities (related to the 't~Hooft--Feynman gauge of the\nquantum fields in the background-field gauge)\nrather than singling out ``its gauge-invariant part'' in any sense.}\nTo facilitate NLO calculations in practice, we have implemented our renormalization schemes\nfor the THDM into a \\textsc{FeynArts} \\cite{Hahn2001} model file, so that amplitudes and squared matrix elements can be generated straightforwardly. \nFinally, we apply the proposed renormalization schemes in the NLO calculation of the partial decay width of the \nlighter CP-even Higgs boson decaying into four fermions, $\\mathswitchr h \\to \\mathswitchr W\\PW\/\\mathswitchr Z\\PZ\\to 4f$,\na process class that is a cornerstone in the experimental determination of Higgs-boson\ncouplings, \n{but for which electroweak corrections in the THDM are not yet known in the literature.}\nThe impact of NLO corrections on Higgs couplings in the THDM was, for instance, investigated more\nglobally in \\citeres{Kanemura:2014dja,Kanemura:2015mxa}.\n{However, a full set of electroweak corrections to all Higgs-boson decay processes in the THDM \ndoes not yet exist in the literature, so that current predictions (see, e.g., \\citere{Harlander:2013qxa})\nfor THDM Higgs analyses globally neglect electroweak higher-order effects.}\n{Our calculation, thus, contributes to overcome this shortcoming;\nelectroweak corrections to some $1\\to2$ particle decays of heavy\nHiggs bosons were presented in \\citere{Krause:2016oke}.}%\n\\footnote{Since we consider the decays of the light Higgs boson $\\mathswitchr h$ via\n$\\mathswitchr W$- or $\\mathswitchr Z$-boson pairs, where at least one of the gauge bosons\nis off its mass shell, we have to consider the full $1\\to4$ process\nwith all off-shell and decay effects, rendering a comparison to\nresults on $\\mathswitchr H\\to\\mathswitchr W\\PW\/\\mathswitchr Z\\PZ$ not meaningful.}\n\nThe structure of the paper is as follows. We introduce the four considered types of THDMs\nand our conventions in Sect.~\\ref{sec:THDM}, and the derivation of the counterterm Lagrangian is performed in Sect.~\\ref{sec:CTLagrangian}. Afterwards we fix the renormalization constants with renormalization conditions (Sect.~\\ref{sec:renconditions}). The on-shell conditions, where the renormalized parameters correspond to measurable quantities, are described and applied in Sect.~\\ref{sec:onshellrenconditions}. The renormalization constants of parameters that do not directly correspond to physical quantities are fixed in the $\\overline{\\mathrm{MS}}$ scheme, so that they contain only the UV divergences and no finite terms. We describe different renormalization schemes based on different definitions of the $\\overline{\\mathrm{MS}}$-renormalized parameters in Sect.~\\ref{sec:diffrenschemes}. \nThe RGEs of the {$\\overline{\\mathrm{MS}}$}{}-renormalized parameters \nare derived and numerically solved in Sect.~\\ref{sec:running}. \nThe implementation of the results into an automated matrix element generator is described in Sect.~\\ref{sec:FAmodelfile}, and numerical results for \nthe partial decay width $\\mathswitchr h\\rightarrow 4 f$ are \npresented in Sect.~\\ref{se:numerics}. \nFinally, we conclude in Sect.~\\ref{se:conclusion}, \nand further details on the renormalization prescription as well as \nsome counterterms are given in the appendix.\n\n\\section{The Two-Higgs-Doublet Model}\n\\label{sec:THDM}\n\nThe Lagrangian of the THDM, $\\mathcal{L}_\\mathrm{THDM}$, is composed of the following parts,\n\\begin{align}\n\\mathcal{L}_\\mathrm{THDM}=\\mathcal{L}_\\mathrm{Gauge}+\\mathcal{L}_\\mathrm{Fermion}+\\mathcal{L}_\\mathrm{Higgs}+\\mathcal{L}_\\mathrm{Yukawa}+\\mathcal{L}_\\mathrm{Fix}+\\mathcal{L}_\\mathrm{Ghost}. \n\\label{eq:LTHDM}\n\\end{align}\nThe gauge, fermionic, gauge-fixing, and ghost parts \ncan be obtained in a straightforward way from the SM ones, e.g., \ngiven in Ref.~\\cite{Denner:1991kt}. The Higgs Lagrangian and the Yukawa couplings to the fermions are discussed in the following and are mostly affected by the additional degrees of freedom \nof the THDM. A very elaborate and complete discussion of the THDM Higgs and Yukawa Lagrangians, including general and specific cases, can, e.g., \nbe found in Refs.~\\cite{Gunion:1989we,Branco:2011iw}.\n\n\\subsection{The Higgs Lagrangian}\n\\label{sec:HiggsL}\n\nThe Higgs Lagrangian, $ \\mathcal{L}_\\mathrm{Higgs}$, contains the kinetic terms and a potential $V$,\n\\begin{align}\n \\mathcal{L}_\\mathrm{Higgs}= (D_\\mu\\Phi_1)^\\dagger (D^\\mu \\Phi_1)+(D_\\mu\\Phi_2)^\\dagger (D^\\mu \\Phi_2)- V(\\Phi_1^\\dagger \\Phi_1,\\Phi_2^\\dagger \\Phi_2,\\Phi_2^\\dagger \\Phi_1,\\Phi_1^\\dagger \\Phi_2),\n\\label{eq:Lhiggs}\n\\end{align}\nwith the complex scalar doublets $\\Phi_{1,2}$ \nof hypercharge $Y_\\mathrm{W}=1$, \n\\begin{align}\n\\Phi_1=\\begin{pmatrix} \\phi_1^+ \\\\ \\phi_1^0\\end{pmatrix}, \\qquad \\Phi_2=\\begin{pmatrix}\\phi_2^+ \\\\ \\phi^0_2\\end{pmatrix},\n\\label{eq:doublets}\n\\end{align}\nand the covariant derivative \n\\begin{align}\n D_\\mu=\\partial_\\mu \\mp \\mathrm{i} g_2 I^a_\\mathrm{W} W^a_\\mu +\\mathrm{i} g_1 \\frac{Y_\\mathrm{W}}{2}B_\\mu,\n\\label{eq:EWcovderivative}\n\\end{align}\nwhere $I^a_\\mathrm{W}$ ($a = 1,2,3$) are the generators of the weak isospin. \nThe SU(2) and \nU(1) gauge fields are denoted \n$W^a_\\mu$ and $B_\\mu$ with the corresponding gauge couplings $g_2$ and $g_1$, respectively.\nThe sign in the \n$g_2$~term is negative in the conventions of B\\\"ohm, Hollik and Spiesberger \n(BHS)~\\cite{Denner:1991kt,Bohm:1986rj} and positive in the convention of Haber and Kane \n(HK)~\\cite{Haber:1984rc}. We implemented both sign conventions, but used the former one as default. \nIn general, the potential involves all hermitian functions of the two doublets up to dimension four and can be parameterized in the most general case as follows \\cite{Wu1994,Branco:2011iw},\n\\begin{align}\nV=&\\,m^2_{11} \\Phi^{\\dagger}_1 \\Phi_1+ m^2_{22} \\Phi^\\dagger_2 \\Phi_2\n-\\left[m^2_{12} \\Phi^\\dagger_1 \\Phi_2 + h.c.\\right]\\nonumber\\\\\n&+ \\frac{1}{2} \\lambda_1 (\\Phi^\\dagger_1 \\Phi_1)^2+\\frac{1}{2} \\lambda_2 (\\Phi^\\dagger_2 \\Phi_2)^2+ \\lambda_3 (\\Phi^\\dagger_1 \\Phi_1)(\\Phi^\\dagger_2 \\Phi_2)+ \\lambda_4 (\\Phi^\\dagger_1 \\Phi_2)(\\Phi^\\dagger_2 \\Phi_1)\\nonumber\\\\\n&+\\left[\\frac{1}{2} \\lambda_5(\\Phi^\\dagger_1 \\Phi_2)^2\n+(\\lambda_6 \\Phi^\\dagger_1\\Phi_1+\\lambda_7 \\Phi^\\dagger_2 \\Phi_2)\\Phi^\\dagger_1 \\Phi_2+h.c.\\right]\\label{eq:THDMpot}.\n\\end{align}\nThe parameters $m_{11}^2,m_{22}^2,\\lambda_1,\\lambda_2,\\lambda_3,\\lambda_4$ are real, while the parameters $m_{12}^2,\\lambda_5,\\lambda_6,\\lambda_7$ are complex, yielding a total number of 14 real degrees of freedom for the potential. However, the component fields of the two Higgs doublets $\\Phi_1$ and $\\Phi_2$ do not correspond to mass eigenstates, and these doublets can be redefined using an SU(2) transformation without changing the physics, so that only 11 physical degrees of freedom remain~\\cite{Branco:2011iw}.\nFor each Higgs doublet we demand that the fields develop a\nvacuum expectation value (vev) in the neutral component,\n\\begin{align}\n \\langle \\Phi_1 \\rangle&=\\langle 0| \\Phi_1|0\\rangle =\\begin{pmatrix}0 \\\\ \\frac{v_1}{\\sqrt{2} } \\end{pmatrix},&\n \\langle \\Phi_2 \\rangle=\\langle 0| \\Phi_2|0\\rangle =\\begin{pmatrix} 0 \\\\ \\frac{v_2}{\\sqrt{2}} \\end{pmatrix}.\n\\end{align}\nIt is non-trivial that such a stable minimum of the potential exists, \nrestricting the allowed parameter space already strongly~\\cite{Ferreira:2004yd}. \nIn general, the vevs are complex (with a significant relative phase).\nThe Higgs doublets can be decomposed as follows,\n\\begin{align}\n\\Phi_1=\\begin{pmatrix} \\phi_1^+ \\\\ \\frac{1}{\\sqrt{2}}(\\eta_1+i \\chi_1+v_1)\\end{pmatrix}, && \\Phi_2=\\begin{pmatrix}\\phi_2^+ \\\\ \\frac{1}{\\sqrt{2}}(\\eta_2+i \\chi_2+v_2)\\end{pmatrix},\n\\label{eq:decom}\n\\end{align}\nwith the charged fields $\\phi_1^+,\\phi_2^+$, the neutral CP-even fields $\\eta_1,\\eta_2$, and the neutral CP-odd fields $\\chi_1,\\chi_2$.\n\n\\paragraph{Additional constraints:}\nSince the 11-dimensional parameter space of the potential is too large for early experimental analyses, we restrict the model in our analysis by imposing two additional conditions, motivated by experimental results:\n\\begin{itemize}\n\\item absence of flavour-changing neutral currents at tree level,\n\\item CP conservation in the Higgs sector (even though this holds only approximately).\n\\end{itemize}\nThe former requirement can be ensured by adding a discrete $\\mathbb{Z}_2$ symmetry $\\Phi_1 \\rightarrow - \\Phi_1$ (see Sect.~\\ref{sec:YukawaL}). This condition implies that the parameters $\\lambda_6$ and $\\lambda_7$ vanish. Permitting operators of dimension two that\nviolate this $\\mathbb{Z}_2$ symmetry softly, non-zero values of $m_{12}$ are still allowed~\\cite{HHGuide:1990, Gunion:2002zf}.\nConcerning the second condition, the potential is CP-conserving if and only if a basis of the Higgs doublets exists in which all parameters and the vevs are real~\\cite{Gunion:2005ja}. For our description we assume that a transformation to such a basis has been done already (if the parameters or vevs were initially complex), so that we only have to deal with real parameters. This renders $m_{12}$ and $\\lambda_5$ real.\nHowever, at higher orders in perturbation theory CP-breaking terms and complex phases in the Higgs sector are generated radiatively through loop contributions\ninvolving the quark mixing matrix. For our NLO analysis, this does not present a problem as they appear only beyond NLO in the specific processes we consider.\nIn addition we assume that a basis of the doublets is chosen in which $v_1,v_2>0$ (which is always possible as a redefinition $\\Phi_i \\to - \\Phi_i$ changes the sign of the vacuum expectation value).\nThe potential~\\eqref{eq:THDMpot} has then the following form,\n\\begin{align}\nV={}&m^2_{11} \\Phi^{\\dagger}_1 \\Phi_1+ m^2_{22} \\Phi^\\dagger_2 \\Phi_2\n-m^2_{12} (\\Phi^\\dagger_1 \\Phi_2 + \\Phi^\\dagger_2 \\Phi_1)\\nonumber\\\\ \n&+ \\frac{1}{2} \\lambda_1 (\\Phi^\\dagger_1 \\Phi_1)^2+\\frac{1}{2} \\lambda_2 (\\Phi^\\dagger_2 \\Phi_2)^2+ \\lambda_3 (\\Phi^\\dagger_1 \\Phi_1)(\\Phi^\\dagger_2 \\Phi_2)+ \\lambda_4 (\\Phi^\\dagger_1 \\Phi_2)(\\Phi^\\dagger_2 \\Phi_1)\\nonumber\\\\ \n&+\\frac{1}{2} \\lambda_5\\left[(\\Phi^\\dagger_1 \\Phi_2)^2+(\\Phi^\\dagger_2 \\Phi_1)^2\\right]. \n\\label{eq:lambdapara}\n\\end{align}\nExpanding the potential using the decomposition~\\eqref{eq:decom} and ordering terms with respect to powers of the fields, leads to the form\n\\begin{align}\nV={}& -t_{\\eta_1} \\eta_1 - t_{\\eta_2} \\eta_2\\nonumber\\\\\n&+\\frac{1}{2} \\,(\\eta_1,\\eta_2) \\mathbf{M}_\\eta \\begin{pmatrix}\\eta_1 \\\\ \\eta_2 \\end{pmatrix}\n+\\frac{1}{2} (\\chi_1,\\chi_2)\\, \\mathbf{M}_\\chi \\begin{pmatrix}\\chi_1 \\\\ \\chi_2 \\end{pmatrix}\n+(\\phi_1^+,\\phi^+_2) \\, \\mathbf{M}_\\phi \\begin{pmatrix}\\phi^-_1 \\\\ \\phi^-_2 \\end{pmatrix}+\\ldots,\\label{eq:potinmatrix}\n\\end{align}\nwith the tadpole terms proportional to the tadpole parameters $t_{\\eta_1}$, $t_{\\eta_2}$ and linear in the fields. The mass terms contain the mass matrices \n$\\mathbf{M}_\\eta$, $\\mathbf{M}_\\chi$, and $ \\mathbf{M}_\\phi$ and are quadratic in the CP-even, CP-odd, and charged Higgs-boson fields, respectively. Terms cubic or quartic in the fields are suppressed in the notation here. Only the neutral CP-even scalar fields can develop non-vanishing tadpole terms, since they carry the quantum numbers of the vacuum. Further, in the mass terms, only particles with the same quantum numbers can mix, so that the three different \ntypes of scalars (neutral CP-even, neutral CP-odd, and charged) do not mix with one another. Of course, through the cubic and quartic terms, which are not shown here, these particles interact with each other. The tadpole parameters are\n\\begin{subequations}\n\\begin{align}\nt_{\\eta_1}=- m^2_{11} v_1- \\lambda_1v_1^3\/2 + v_2 (m^2_{12}-\\lambda_{345}v_1 v_2\/2),\\\\\nt_{\\eta_2}=- m^2_{22} v_2- \\lambda_2v_2^3\/2 + v_1 (m^2_{12}-\\lambda_{345} v_1 v_2\/2),\n\\end{align}%\n\\label{eq:tadpoleterms}%\n\\end{subequations}%\nwhere we introduced the abbreviations $\\lambda_{ij\\ldots}=\\lambda_i+\\lambda_j+\\ldots$, and the mass matrices are given by\n\\begin{subequations}\n\\begin{align}\n\\mathbf{M}_\\eta&=\\begin{pmatrix} m^2_{11}+ 3 \\lambda_1 v_1^2\/2+ \\lambda_{345} v_2^2\/2 & -m^2_{12}+ \\lambda_{345} v_1 v_2\\\\ -m^2_{12}+ \\lambda_{345} v_1 v_2& m^2_{22}+ 3 \\lambda_2 v_2^2\/2+ \\lambda_{345} v_1^2\/2\\end{pmatrix},\n\\\\\n\\mathbf{M}_\\chi&=\\begin{pmatrix} m^2_{11}+ \\lambda_1 v_1^2\/2+ (\\lambda_{34}-\\lambda_{5}) v_2^2\/2 & -m^2_{12}+ \\lambda_{5} v_1 v_2\\\\ \n\\hspace{-0pt}-m^2_{12}+ \\lambda_{5} v_1 v_2& \n\\hspace{-0pt}m^2_{22}+ \\lambda_2 v_2^2\/2+ (\\lambda_{34}-\\lambda_{5}) v_1^2\/2\\end{pmatrix},\n\\\\\n\\mathbf{M}_\\phi&=\n\\begin{pmatrix} m^2_{11}+ \\lambda_1 v_1^2\/2+ \\lambda_{3} v_2^2\/2 \n& -m^2_{12}+ \\lambda_{45} v_1 v_2\/2 \\\\ \n-m^2_{12}+ \\lambda_{45} v_1 v_2\/2 & m^2_{22}+ \\lambda_2 v_2^2\/2+ \\lambda_{3} v_1^2\/2\n\\end{pmatrix}.\n\\label{eq:massmatricesfirst}\n\\end{align}\n\\end{subequations}\nThe fields can be transformed into their mass eigenstate basis via\n\\begin{subequations}\n\\label{eq:rotations}\n\\begin{align}\n\\label{eq:Hhrot}\n\\begin{pmatrix}\\eta_1 \\\\ \\eta_2\\end{pmatrix} &\n= \\begin{pmatrix}\\cos{\\alpha} & -\\sin{\\alpha}\\\\ \\sin{\\alpha}& \\cos{\\alpha}\\end{pmatrix} \\begin{pmatrix}H \\\\ h\\end{pmatrix},\n\\\\\n\\begin{pmatrix}\\chi_1 \\\\ \\chi_2\\end{pmatrix}&\n= \\begin{pmatrix}\\cos{\\beta_n} & -\\sin{\\beta_n}\\\\ \\sin{\\beta_n}& \\cos{\\beta_n}\\end{pmatrix} \\begin{pmatrix}G_0 \\\\ A_0\\end{pmatrix},\n\\\\\n\\begin{pmatrix}\\phi^\\pm_1 \\\\ \\phi^\\pm_2\\end{pmatrix}&\n= \\begin{pmatrix}\\cos{\\beta_c} & -\\sin{\\beta_c}\\\\ \\sin{\\beta_c}& \\cos{\\beta_c}\\end{pmatrix} \\begin{pmatrix}G^\\pm \\\\ H^\\pm\\end{pmatrix}.\n\\end{align}\n\\end{subequations}\nwhere $h$, $H$ correspond to the CP-even, $A_0$ to the CP-odd, and $H^\\pm$ to the charged mass eigenstates\\footnote{In order to avoid a conflict in our notation, we define $\\alpha_\\mathrm{em}=e^2\/(4\\pi)$ as electromagnetic coupling constant and consistently keep the symbol $\\alpha$ for the rotation angle.}. The fields $G^\\pm$, \n$G_0$ correspond to the Goldstone bosons.\nAfter a rotation of the fields, the potential has the following form,\n\\begin{align}\n\\label{eq:NLObarepot}\nV=& -t_\\mathswitchr H H - t_\\mathswitchr h h\\nonumber\\\\*\n&+\\frac{1}{2} (H,h) \\begin{pmatrix}\\mathswitch {M_\\PH}^2 & M^2_{\\mathswitchr H\\mathswitchr h}\\\\ M^2_{\\mathswitchr H\\mathswitchr h}& \\mathswitch {M_\\Ph}^2\\end{pmatrix} \\begin{pmatrix}H \\\\ h \\end{pmatrix}\n+\\frac{1}{2} (G_0,A_0) \\begin{pmatrix}M^2_\\mathswitchr {G_0} & M^2_{\\mathswitchr {G_0}\\mathswitchr {A_0}}\\\\ M^2_{\\mathswitchr {G_0}\\mathswitchr {A_0}}& M^2_{\\mathswitchr {A_0}}\\end{pmatrix} \\begin{pmatrix}G_0 \\\\ A_0 \\end{pmatrix}\\nonumber\\\\\n&+ (G^+,H^+) \\begin{pmatrix}M_\\mathswitchr {G^+}^2 & M^2_{\\mathswitchr {G}\\mathswitchr {H^+}}\\\\ M^2_{\\mathswitchr {G}\\mathswitchr {H^+}}& \\mathswitch {M_\\PHP}^2\\end{pmatrix} \\begin{pmatrix}G^- \\\\ H^- \\end{pmatrix}+ \\textrm{interaction terms} \n\\end{align}\nwith the \ntadpole parameters\n\\begin{align}\n\\label{eq:tadpoledefs}\nt_\\mathswitchr H=\\mathswitch {c_\\alpha} t_{\\eta_1}+\\mathswitch {s_\\alpha} t_{\\eta_2}, \\qquad\nt_\\mathswitchr h=-\\mathswitch {s_\\alpha} t_{\\eta_1}+\\mathswitch {c_\\alpha} t_{\\eta_2},\n\\end{align}\nwhere general abbreviations for the trigonometric functions $s_x\\equiv\\sin x $, $c_x\\equiv\\cos x$, $t_x\\equiv\\tan x$ are introduced.\nAfter the elimination of $m_{11}$, $m_{22}$ using the above equations, the entries of the mass matrices contain also the tadpole parameters $t_h$, $t_H$. \nUsing \n\\begin{align}\nv^2=v_1^2+v_2^2, \\qquad \n\\tan{\\beta} = \\frac{v_2}{v_1},\n\\label{eq:tanbeta}\n\\end{align}\nwe obtain for the mass parameters of the CP-even Higgs bosons,\n\\begin{subequations}\n\\label{eq:neutralmassdef}\n\\begin{align}\n\\mathswitch {M_\\PH}^2={}&\\frac{2s_{\\alpha-\\beta}^2}{s_{2\\beta}} m^2_{12}\n+ \\frac{v^2}{2} \\left( 2\\lambda_1 c_\\beta^2 c_\\alpha^2 +2\\lambda_2 s_\\alpha^2 s_\\beta^2 +s_{2 \\alpha} s_{2\\beta} \\lambda_{345}\\right)\n-2t_\\mathswitchr H \\frac{s_\\alpha^3 c_\\beta+c_\\alpha^3 s_\\beta }{ v s_{2 \\beta }}\n-t_\\mathswitchr h\\frac{s_{2 \\alpha } s_{\\alpha -\\beta }}{v s_{2 \\beta }},\n\\\\\n\\mathswitch {M_\\Ph}^2={}&\\frac{2c_{\\alpha - \\beta}^2}{s_{2 \\beta}} m^2_{12}\n+ \\frac{v^2}{2} \\left( 2\\lambda_1 c_\\beta^2 s_\\alpha^2 +2\\lambda_2 c_\\alpha^2 s_\\beta^2-s_{2 \\alpha} s_{2\\beta} \\lambda_{345}\\right)\n-t_\\mathswitchr H\\frac{s_{2 \\alpha } c_{\\alpha -\\beta }}{v s_{2 \\beta }}\n-2t_\\mathswitchr h\\frac{c_\\alpha^3 c_\\beta-s_\\alpha^3 s_\\beta }{ v s_{2\\beta }},\n\\\\\nM_{\\mathswitchr H\\mathswitchr h}^2={}& \\frac{s_{2(\\alpha-\\beta)}}{s_{2 \\beta}}m^2_{12}\n+ \\frac{v^2}{2} \\left[ s_{2\\alpha} (-c_\\beta^2 \\lambda_1+ s_\\beta^2 \\lambda_2)+ s_{2 \\beta} c_{2\\alpha} \\lambda_{345}\\right]\n-t_\\mathswitchr H\\frac{s_{2 \\alpha } s_{\\alpha -\\beta }}{v s_{2 \\beta }}\n-t_\\mathswitchr h\\frac{s_{2 \\alpha } c_{\\alpha -\\beta }}{v s_{2 \\beta }},\n\\label{eq:Hhmassmixing}\n\\end{align}\n\\end{subequations}\nthe ones of the mass matrix of the CP-odd Higgs fields result in\n\\begin{subequations}\n\\label{eq:cpmassdef}\n\\begin{align}\n\\mathswitch {M_\\PAO}^2&=2c_{\\beta-\\beta_n}^2\\Big(\\frac{m^2_{12}}{s_{2\\beta}}- \\frac{\\lambda_5 v^2}{2}\\Big) \n- 2t_\\mathswitchr H \\frac{c_{\\beta_n}^2 c_\\beta s_\\alpha + s_{\\beta_n}^2 s_\\beta c_\\alpha}{v s_{2\\beta}}\n- 2t_\\mathswitchr h \\frac{c_{\\beta_n}^2 c_\\beta c_\\alpha - s_{\\beta_n}^2 s_\\beta s_\\alpha}{v s_{2\\beta}}, \n\\label{eq:A0massdef}\\\\\nM_{\\mathswitchr {G_0}}^2&=2s_{\\beta-\\beta_n}^2\\Big(\\frac{m^2_{12}}{s_{2\\beta}}- \\frac{\\lambda_5 v^2}{2}\\Big)\n- 2t_\\mathswitchr H \\frac{s_{\\beta_n}^2 c_\\beta s_\\alpha + c_{\\beta_n}^2 s_\\beta c_\\alpha}{v s_{2\\beta}}\n- 2t_\\mathswitchr h \\frac{s_{\\beta_n}^2 c_\\beta c_\\alpha - c_{\\beta_n}^2 s_\\beta s_\\alpha}{v s_{2\\beta}},\\\\\nM_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2&=-s_{2(\\beta-\\beta_n)}\\Big(\\frac{m^2_{12}}{s_{2\\beta}}- \\frac{\\lambda_5 v^2}{2}\\Big)- t_\\mathswitchr H \\frac{s_{2\\beta_n} s_{\\alpha-\\beta}}{v s_{2\\beta}}- t_\\mathswitchr h \\frac{s_{2\\beta_n} c_{\\alpha-\\beta}}{v s_{2\\beta}}, \\label{eq:goldstonemassmixing}\n\\end{align}\n\\end{subequations}\nand the ones of the mass matrix of the charged Higgs-boson fields are\n\\begin{subequations}\n\\label{eq:chargedmassdef}\n\\begin{align}\n\\mathswitch {M_\\PHP}^2&=c_{\\beta-\\beta_c}^2\\Big[\\frac{2m^2_{12}}{s_{2\\beta}}-\\frac{v^2}{2}(\\lambda_4+\\lambda_5) \\Big]\n- 2t_\\mathswitchr H \\frac{c_{\\beta_c}^2 c_\\beta s_\\alpha + s_{\\beta_c}^2 s_\\beta c_\\alpha}{v s_{2\\beta}}\n- 2t_\\mathswitchr h \\frac{c_{\\beta_c}^2 c_\\beta c_\\alpha - s_{\\beta_c}^2 s_\\beta s_\\alpha}{v s_{2\\beta}},\n\\label{eq:Hpmmassdef}\\\\\nM_\\mathswitchr {G^+}^2&=s_{\\beta-\\beta_c}^2\\Big[\\frac{2m^2_{12}}{s_{2\\beta}}-\\frac{v^2}{2}(\\lambda_4+\\lambda_5) \\Big]\n- 2t_\\mathswitchr H \\frac{s_{\\beta_c}^2 c_\\beta s_\\alpha + c_{\\beta_c}^2 s_\\beta c_\\alpha}{v s_{2\\beta}}\n- 2t_\\mathswitchr h \\frac{s_{\\beta_c}^2 c_\\beta c_\\alpha - c_{\\beta_c}^2 s_\\beta s_\\alpha}{v s_{2\\beta}},\n\\\\\nM_{\\mathswitchr {G}\\mathswitchr {H^+}}^2&=-\\frac{s_{2(\\beta-\\beta_c)}}{2} \\Big[\\frac{2m^2_{12}}{s_{2\\beta}}-\\frac{v^2}{2}(\\lambda_4+\\lambda_5) \\Big]\n- t_\\mathswitchr H \\frac{s_{2\\beta_c} s_{\\alpha-\\beta}}{v s_{2 \\beta}}\n- t_\\mathswitchr h \\frac{s_{2\\beta_c} c_{\\alpha-\\beta}}{v s_{2 \\beta}}.\n\\label{eq:chargedgoldstonemassmixing}\n\\end{align}\n\\end{subequations}\nAt tree level we demand vanishing tadpole terms corresponding to $t_\\mathswitchr H=t_\\mathswitchr h=0$ and diagonal propagators, so that the mixing terms proportional to $M_{\\mathrm{Hh}}^2$, $M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2$, $M_{\\mathswitchr {G}\\mathswitchr {H^+}}^2$ vanish at LO. This yields $\\beta_c=\\beta_n=\\beta$ and fixes the angle $\\alpha$ as well.\nAt NLO such a diagonalization is not possible, as the propagators receive also mixing contributions from the field renormalization and from \n(momentum-dependent) one-loop diagrams, so that there is no distinct condition to define the mixing angles. Therefore we keep the bare mass mixing parameters $M_{\\mathrm{Hh}}^2$, $M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2$, $M_{\\mathswitchr {G}\\mathswitchr {H^+}}^2$ and the tadpole terms $t_\\mathswitchr H$, $t_\\mathswitchr h$ in this section, and \nspecify defining conditions for the bare parameters $\\alpha$, $\\beta_n$, $\\beta_c$ and the tadpole terms later.\nWith the above equations, $m_{12}$, $\\lambda_1$, $\\lambda_2$, $\\lambda_3$, $\\lambda_4$ can be traded for the masses of the physical bosons $M_\\mathswitchr H$, $M_\\mathswitchr h$, $M_\\mathswitchr {A_0}$, $\\mathswitch {M_\\PHP}$, and the mixing angle $\\alpha$. The parameter $\\lambda_5$ cannot be replaced by a mass or a mixing angle as it appears only in cubic and quartic Higgs couplings and acts like an additional coupling constant. Explicit relations can be obtained by inverting Eqs.~\\eqref{eq:neutralmassdef}, \n\\eqref{eq:A0massdef}, and~\\eqref{eq:Hpmmassdef}.\nThe other parameters are related to the masses by\n\\begin{subequations}\n\\label{eq:repara}\n\\begin{align}\n\\lambda_1={}&\\frac{1}{\\mathswitch {c_\\beta}^2 v^2} \\left[\\mathswitch {c_\\alpha}^2 \\mathswitch {M_\\PH}^2+\\mathswitch {s_\\alpha}^2 \\mathswitch {M_\\Ph}^2- M_{\\mathswitchr H\\mathswitchr h}^2 s_{2\\alpha}-\\sb^2 \\biggl(\\frac{\\mathswitch {M_\\PAO}^2}{ c_{\\beta-\\beta_n}^2}+\\lambda_5 v^2\\biggr) \\right]\n\\nonumber\\\\\n&+t_\\mathswitchr H\\frac{c_{{\\beta_{n}}} \\left(2 s_{\\beta } s_{{\\beta_{n}}} c_{\\alpha }+c_{\\alpha +\\beta }c_{\\beta_n}\\right)}{v^3 c_{\\beta }^2 c_{\\beta -{\\beta_{n}}}^2}\n -t_\\mathswitchr h \\frac{c_{\\beta_{n}} \\left(2 s_{\\beta } s_{\\beta_{n}} s_{\\alpha }+s_{\\alpha +\\beta }c_{\\beta_n}\\right)}{v^3 c_{\\beta }^2 c_{\\beta -\\beta_{n}}^2},\n\\\\\n\\lambda_2={}&\\frac{1}{\\sb^2 v^2} \\left[\\mathswitch {c_\\alpha}^2 \\mathswitch {M_\\Ph}^2+\\mathswitch {s_\\alpha}^2 \\mathswitch {M_\\PH}^2+ M_{\\mathswitchr H\\mathswitchr h}^2 s_{2\\alpha}-\\mathswitch {c_\\beta}^2 \\biggl(\\frac{\\mathswitch {M_\\PAO}^2}{c_{\\beta-\\beta_n}^2}+\\lambda_5 v^2\\biggr) \\right]\\nonumber\\\\\n&+ t_\\mathswitchr H \\frac{s_{\\beta_n} (2 c_\\beta c_{\\beta_n} s_\\alpha- c_{\\alpha+\\beta} s_{\\beta_n})}{v^3 s_\\beta^2 c_{\\beta-\\beta_n}^2}+ t_\\mathswitchr h \\frac{s_{\\beta_n} (2 c_\\beta c_{\\beta_n} c_\\alpha+ s_{\\alpha+\\beta} s_{\\beta_n})}{v^3 s_\\beta^2 c_{\\beta-\\beta_n}^2},\n\\\\\n\\lambda_3={}&\\frac{1}{v^2 s_{2 \\beta}} \\left[s_{2\\alpha}(\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2) + 2 c_{2\\alpha}M_{\\mathswitchr H\\mathswitchr h}^2\\right]-\\frac{\\mathswitch {M_\\PAO}^2}{v^2c_{\\beta-\\beta_n}^2}+\\frac{2 \\mathswitch {M_\\PHP}^2}{v^2c_{\\beta-\\beta_c}^2}-\\lambda_5\\nonumber \n\\\\\n&+\\frac{2t_\\mathswitchr H}{v^3 s_{2\\beta }} \n\\left[s_{\\beta } c_{\\alpha } \\bigg(\\frac{2 s^2_{\\beta_c}}{c_{\\beta -\\beta_c}^2}-\\frac{s^2_{\\beta_n}}{c_{\\beta -\\beta_n}^2}\\bigg)+c_{\\beta } s_{\\alpha }\n \\bigg(\\frac{2 c^2_{\\beta_c}}{c_{\\beta -\\beta_c}^2}-\\frac{c^2_{\\beta_n}}{c_{\\beta -\\beta_n}^2}\\bigg)\\right]\n\\nonumber\\\\\n&+ \\frac{2t_\\mathswitchr h}{v^3 s_{2\\beta }} \n\\left[c_{\\beta } c_{\\alpha }\\bigg(\\frac{2 c^2_{\\beta_c}}{c_{\\beta -\\beta_c}^2}-\\frac{c^2_{\\beta_n}}{c_{\\beta -\\beta_n}^2}\\bigg)+s_{\\beta } s_{\\alpha }\n \\bigg(\\frac{s^2_{\\beta_n}}{c_{\\beta -\\beta_n}^2}-\\frac{2 s^2_{\\beta_c}}{c_{\\beta -\\beta_c}^2}\\bigg)\\right]\n\\label{eq:Lambda3NLO}\\\\\n\\lambda_4={}&\\lambda_5+\\frac{2 \\mathswitch {M_\\PAO}^2}{v^2 c_{\\beta -\\beta _n}^2}-\\frac{2 \\mathswitch {M_\\PHP}^2}{v^2 c_{\\beta -\\beta _c}^2}+\\frac{2 t_\\mathrm{H} s_{\\beta _c-\\beta _n} \\left(s_{\\alpha +\\beta -\\beta _c-\\beta\n _n}-s_{\\beta -\\alpha } c_{\\beta _c-\\beta _n}\\right)}{v^3 c_{\\beta -\\beta _c}^2 c_{\\beta -\\beta _n}^2}\\nonumber\\\\\n&+\\frac{2 t_\\mathrm{h} s_{\\beta _c-\\beta _n} \\left(c_{\\alpha +\\beta -\\beta _c-\\beta\n _n}+c_{\\beta -\\alpha } c_{\\beta _c-\\beta _n}\\right)}{v^3 c_{\\beta -\\beta _c}^2 c_{\\beta -\\beta _n}^2},\\\\\nm^2_{12}={}&\\frac{1}{2}\\lambda_5 v^2 s_{2\\beta }\n+\\frac{\\mathswitch {M_\\PAO}^2 s_{2\\beta }}{2c_{\\beta -\\beta _n}^2}\n+\\frac{t_\\mathrm{H} \\left(s_{\\beta } c_{\\alpha } s_{\\beta _n}^2+c_{\\beta } s_{\\alpha } c_{\\beta_n}^2\\right)}{v c_{\\beta -\\beta _n}^2}\n+\\frac{t_\\mathrm{h} \\left(c_{\\beta } c_{\\alpha } c_{\\beta _n}^2-s_{\\beta }\n s_{\\alpha } s_{\\beta _n}^2\\right)}{v c_{\\beta -\\beta _n}^2}.\\label{eq:m12NLO}\n\\end{align}\n\\end{subequations}\nThe tree-level relations are\neasily obtained by setting $\\beta_c=\\beta_n=\\beta$ and $t_\\mathswitchr H=t_\\mathswitchr h=M_{\\mathrm{Hh}}^2=M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2=M_{\\mathswitchr {G}\\mathswitchr {H^+}}^2=0$,\n\\begin{subequations}\n\\begin{align}\n\\lambda_1&={}\\frac{1}{\\mathswitch {c_\\beta}^2 v^2} \\left[\\mathswitch {c_\\alpha}^2 \\mathswitch {M_\\PH}^2+\\mathswitch {s_\\alpha}^2\n\\mathswitch {M_\\Ph}^2-\\sb^2 (\\mathswitch {M_\\PAO}^2+\\lambda_5 v^2) \\right],\\\\\n\\lambda_2&={}\\frac{1}{\\sb^2 v^2} \\left[\\mathswitch {s_\\alpha}^2 \\mathswitch {M_\\PH}^2+\\mathswitch {c_\\alpha}^2\n\\mathswitch {M_\\Ph}^2-\\mathswitch {c_\\beta}^2 (\\mathswitch {M_\\PAO}^2+\\lambda_5 v^2) \\right],\\\\\n\\lambda_3&={}\\frac{s_{2\\alpha}}{s_{2\\beta} v^2}\n(\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2)-\\frac{1}{v^2} (\\mathswitch {M_\\PAO}^2-2\\mathswitch {M_\\PHP}^2)-\\lambda_5,\n\\label{eq:lambda3tree}\\\\\n\\lambda_4&={}\\frac{2(\\mathswitch {M_\\PAO}^2-\\mathswitch {M_\\PHP}^2)}{v^2}+\\lambda_5,\\\\\nm^2_{12}&={}\\frac{s_{2\\beta}}{2} (\\mathswitch {M_\\PAO}^2+\\lambda_5 v^2).\n\\end{align}\n\\end{subequations}\n\n\\paragraph{Parameters of the gauge sector:}\nMass terms of the gauge bosons arise through the interaction of the gauge bosons with the vevs, \nanalogous to the SM. After a rotation into fields corresponding to mass eigenstates, one obtains relations similar to the SM ones:\n\\begin{align}\n\\mathswitch {M_\\PW}=g_2 \\frac{v}{2},&& \\mathswitch {M_\\PZ}=\\frac{v}{2} \\sqrt{g_1^2+g_2^2},&& e= \\frac{g_1g_2}{\\sqrt{g_1^2+g_2^2}},\n\\label{eq:EWmassesTHDM}\n\\end{align}\nwhere the electric unit charge $e$ is identified with the coupling constant of the photon field $A_\\mu$ in the covariant derivative.\nInverting these relations and introducing the weak mixing angle $\\theta_W$ via $\\cos \\theta_W=g_2\/\\sqrt{g_1^2+g_2^2}$, one can replace $v$ and the gauge couplings $g_1$ and $g_2$ by\n\\begin{align}\nv&= \\frac{2 \\mathswitch {M_\\PW} \\mathswitch {s_{\\scrs\\PW}}}{e},& g_1&=\\frac{e}{\\mathswitch {c_{\\scrs\\PW}}},& g_2&=\\frac{e}{\\mathswitch {s_{\\scrs\\PW}}}.\n\\label{eq:EWtrafo}\n\\end{align}\n\n\\paragraph{Mass parameterization:}\nThe relations~\\eqref{eq:tadpoledefs}, \n\\refeq{eq:tanbeta},\n(\\ref{eq:repara}a,b,d,e), and \\eqref{eq:EWtrafo} between the masses, angles, and the basic parameters can be used to reparameterize the Higgs Lagrangian and change the parameters\n\\begin{align}\n\\{p_{\\mathrm{basic}}\\}=\\{\\lambda_1,\\ldots,\\lambda_5,m^2_{11},m^2_{22},m^2_{12},v_1,v_2,g_1,g_2\\},\\label{eq:basicpara}\n\\end{align}\nin favour of the bare mass parameters including one parameter from scalar self-interactions\nwhich we take as $\\lambda_3$, \n\\begin{align}\n\\{p'_{\\mathrm{mass}}\\}=\\{\\mathswitch {M_\\PH},\\mathswitch {M_\\Ph},\\mathswitch {M_\\PAO},\\mathswitch {M_\\PHP},M_\\mathrm{W},\\mathswitch {M_\\PZ},e,\\lambda_5,\\lambda_3,\\beta, t_\\mathrm{H} ,t_\\mathrm{h}\\}\n\\label{eq:physparaNLOprime}\n\\end{align}\nAdditionally, one has to keep in mind that we keep the mixing parameters $\\alpha$, $\\beta_n$, and $\\beta_c$ generic, and they have to be fixed by additional conditions (which will be given later).\nOne can use Eq.~(\\ref{eq:repara}c) to trade $\\lambda_3$ for $\\alpha$, in which case the mixing angle becomes a free parameter of the theory. Then, one obtains the parameter set\n\\begin{align}\n\\{p_{\\mathrm{mass}}\\}=\\{\\mathswitch {M_\\PH},\\mathswitch {M_\\Ph},\\mathswitch {M_\\PAO},\\mathswitch {M_\\PHP},M_\\mathrm{W},\\mathswitch {M_\\PZ},e,\\lambda_5,\\alpha,\\beta, t_\\mathrm{H} ,t_\\mathrm{h} \\}\n\\label{eq:physparaNLO}\n\\end{align}\n\n\\subsection{Yukawa couplings}\n\\label{sec:YukawaL}\n\nThe Higgs mechanism does not only give rise to the gauge-boson mass terms (which are determined by the vevs), \nbut via Yukawa couplings, it introduces\nmasses to chiral fermions. Since both Higgs doublets can couple to fermions, the general Yukawa couplings have the form\n\\begin{align}\n\\mathcal{L}_\\mathrm{Yukawa}=&-\\sum_{k=1,2}\\sum_{i,j} \\left(\\bar{L}'^{\\mathrm{L}}_i \\zeta^{l,k}_{ij} l'^{\\mathrm{R}}_j \\Phi_k+\\bar{Q}'^{\\mathrm{L}}_i \\zeta^{u,k}_{ij} u'^{\\mathrm{R}}_j \\tilde{\\Phi}_k+\\bar{Q}'^{\\mathrm{L}}_i \\zeta^{d,k}_{ij} d'^{\\mathrm{R}}_j \\Phi_k+h.c.\\right),\n\\label{LYuk}\n\\end{align}\nwith the mixing matrices $\\zeta^{f,k}$, $k=1,2$, in generation space for the \ngauge-invariant interactions with $\\Phi_1$ and $\\Phi_2$, respectively, and the generation indices $i, j = 1,2,3$.\nThe left-handed SU(2) doublets of quarks and leptons are denoted \n${Q}'^{\\mathrm{L}}=\\left(u'^{\\mathrm{L}},d'^{\\mathrm{L}}\\right)^{\\mathrm{T}}$ and \n${L}'^{\\mathrm{L}}=\\left(\\nu'^{\\mathrm{L}},l'^{\\mathrm{L}}\\right)^{\\mathrm{T}}$, \nwhile the right-handed up-type quark, down-type quark, and lepton singlets are $u'^{\\mathrm{R}}$, $d'^{\\mathrm{R}}$, and $l'^{\\mathrm{R}}$, respectively. \nThe primes indicate that we deal with fields \nin the interaction basis here;\nfields without primes correspond to mass eigenstates.\nThe field $\\tilde{\\Phi}_k$, $k = 1, 2$, is the charge-conjugated field of $\\Phi_k$.\nSince, in the general THDM, there are two mass mixing matrices for each type $f$ of fermions, flavour-changing neutral currents (FCNC) can occur at tree level, which, however, are experimentally known to be strongly suppressed.\nAccording to the Paschos--Glashow--Weinberg theorem~\\cite{Glashow1977,Paschos1977a}, \nFCNC are absent at tree level\nif each type of fermion couples only to one of the Higgs doublets.\n\n This can be achieved by imposing an additional discrete $\\mathbb{Z}_2$ symmetry. It should be noted that the soft-$\\mathbb{Z}_2$-breaking term proportional to $m_{12}$ in the Higgs potential does not introduce FCNC. \nThe Yukawa Lagrangian reduces then to\n\\begin{align}\n\\mathcal{L}_\\mathrm{Yukawa}=-\\sum_{i,j} \\left(\\bar{L}'^{\\mathrm{L}}_i \\zeta^{l}_{ij} l'^{\\mathrm{R}}_j \\Phi_{n_1}+\\bar{Q}'^{\\mathrm{L}}_i \\zeta^{u}_{ij} u'^{\\mathrm{R}}_j \\tilde{\\Phi}_{n_2}+\\bar{Q}'^{\\mathrm{L}}_i \\zeta^{d}_{ij} d'^{\\mathrm{R}}_j \\Phi_{n_3}+h.c.\\right),\n\\end{align}\nwith $n_i$ being either $1$ or $2$. \nDepending on the exact form of the symmetry,\none distinguishes four types of THDMs. In Type~I models, all fermions couple to one Higgs doublet (conventionally $\\Phi_2$, but this is equivalent to $\\Phi_1$ due to possible basis changes) which can be ensured by demanding a $\\Phi_1 \\rightarrow -\\Phi_1$ symmetry. In Type~II models, down-type fermions couple to the other doublet, which can be enforced by the symmetry $\\Phi_1 \\rightarrow -\\Phi_1$, $d'^{\\mathrm{R}}_j \\to- d'^{\\mathrm{R}}_j$, $l'^{\\mathrm{R}} \\to -l'^{\\mathrm{R}}$. The other two possibilities are called ``lepton-specific'' (Type~X) and ``flipped'' (Type~Y) models. An overview over the couplings and symmetries of the different models is given in Tab.~\\ref{tab:yukcoupl}.\n\\begin{table}\n \\centering\n\\begin{tabular}{|c|c|c|c|c|}\\hline\n&\\hspace{20pt} $u_i$\\hspace{20pt} &\\hspace{20pt} $d_i$\\hspace{20pt} &\\hspace{20pt} $e_i$\\hspace{20pt} & $\\mathbb{Z}_2$ symmetry \\\\\\hline\nType I & $\\Phi_2$ & $\\Phi_2$ & $\\Phi_2$ & $\\Phi_1 \\rightarrow - \\Phi_1$ \\\\\nType II& $\\Phi_2$ & $\\Phi_1$ & $\\Phi_1$ &$(\\Phi_1, d_i,e_i) \\rightarrow -(\\Phi_1, d_i,e_i)$\\\\\nLepton-specific & $\\Phi_2$ & $\\Phi_2$ & $\\Phi_1$ &$(\\Phi_1, e_i) \\rightarrow -(\\Phi_1,e_i)$ \\\\\nFlipped & $\\Phi_2$ & $\\Phi_1$ & $\\Phi_2$ & $(\\Phi_1, d_i) \\rightarrow -(\\Phi_1, d_i)$\\\\\\hline\n\\end{tabular} \n \\caption{Different types of the THDM having in common that only one of the Higgs doublet couples to each type of fermions. This can be achieved by imposing appropriate $\\mathbb{Z}_2$ symmetry charges to the fields. }\n\\label{tab:yukcoupl}\n\\end{table}\nFor each of the fermion types, a redefinition of the fields can be performed in order to get diagonal mass matrices, analogously to the SM case. Similar to the SM, \nthe coupling of fermions to the Z boson is flavour conserving, and a CKM matrix appears in the coupling to the charged gauge bosons. \nBy specifying the model type, the Higgs--fermion interaction is determined, and one can write them, \nwidely following the notation of \nRef.~\\cite{Aoki2009}, as\n\\begin{align}\n \\mathcal{L}_\\mathrm{Yukawa,int}=\n&- \\sum_i \\sum_{f=u,d,l} \\frac{m_{f,i}}{v}\\left( \\xi^f_\\mathswitchr h\\, \\bar{f_i}f_i h+ \\xi^f_\\mathswitchr H \\, \\bar{f}_if_iH\n-2\\mathrm{i} I^3_{\\mathswitchr W,f} \\xi^f_\\mathswitchr {A_0} \\, \\bar{f}_i\\gamma_5 f_i A_0\n- 2\\mathrm{i} I^3_{\\mathswitchr W,f} \\bar{f}_i\\gamma_5 f_i G_0\\right)\\nonumber\\\\\n&-\\sum_{i,j}\\biggl[\\frac{\\sqrt{2}V_{ij}}{v} \\bar{u}_i (-m_{u,i} \\xi^u_\\mathswitchr {A_0} \\omega_- +m_{d,j} \\xi^d_\\mathswitchr {A_0} \\omega_+)\\, d_j H^+ + h.c. \\biggr]\n\\nonumber\\\\\n&- \\sum_i \\biggl[ \\frac{\\sqrt{2}m_{l,i} \\xi^l_\\mathswitchr {A_0}}{v}\\, \\bar{\\nu}_i^{\\mathrm{L}} l_i^{\\mathrm{R}} H^+ + h.c. \\biggr]\n\\nonumber\\\\\n&-\\sum_{i,j}\\biggl[ \\frac{\\sqrt{2}V_{ij}}{v} \\bar{u}_i (-m_{u,i} \\omega_- +m_{d,j} \\omega_+)\\, d_j G^+ + h.c. \\biggr]\n- \\sum_i \\biggl[ \\frac{\\sqrt{2}m_{l,i} }{v}\\, \\bar{\\nu}_i^{\\mathrm{L}} l_i^{\\mathrm{R}} G^+ + h.c. \\biggr],\n\\label{eq:ffH}\n\\end{align}\nwhere $m_{l,i}$, $m_{u,i}$, and $m_{d,i}$ are the lepton, the up-type, and the down-type quark masses, respectively,\nand $V_{ij}$ are the coefficients of the CKM matrix.\nLeft- and right-handed fermion fields, $f^{\\mathrm{L}}$ and $f^{\\mathrm{R}}$, are obtained from the \ncorresponding Dirac spinor $f$ by applying the chirality projectors\n$\\omega_\\pm= ( 1\\pm \\gamma_5)\/2$, i.e.\\ $f=(\\omega_+ +\\omega_-)f=f^{\\mathrm{L}}+f^{\\mathrm{R}}$.\nThe coupling coefficients $\\xi^f_{\\mathswitchr H,\\mathswitchr h,\\mathswitchr {A_0}}$ are defined as the couplings relative to the canonical SM value of $m_f\/v$ and are shown in Tab.~\\ref{tab:yukint}. \n\\begin{table}\n \\centering\n\\begin{tabular}{|c|c|c|c|c|}\\hline\n&Type I & Type II & Lepton-specific & Flipped\\\\\\hline\n$\\xi^l_\\mathswitchr H$ & $\\sin{\\alpha}\/\\sin{\\beta}$ & $\\cos{\\alpha}\/\\cos{\\beta}$ & $\\cos{\\alpha}\/\\cos{\\beta}$ &$\\sin{\\alpha}\/\\sin{\\beta}$\\\\\n$\\xi^u_\\mathswitchr H$ & $\\sin{\\alpha}\/\\sin{\\beta}$ & $\\sin{\\alpha}\/\\sin{\\beta}$ & $\\sin{\\alpha}\/\\sin{\\beta}$ &$\\sin{\\alpha}\/\\sin{\\beta}$ \\\\\n$\\xi^d_\\mathswitchr H$& $\\sin{\\alpha}\/\\sin{\\beta}$ & $\\cos{\\alpha}\/\\cos{\\beta}$ & $\\sin{\\alpha}\/\\sin{\\beta}$ & $\\cos{\\alpha}\/\\cos{\\beta}$ \\\\\\hline\n$\\xi^l_\\mathswitchr h$ & $\\cos{\\alpha}\/\\sin{\\beta}$ & $-\\sin{\\alpha}\/\\cos{\\beta}$ & $-\\sin{\\alpha}\/\\cos{\\beta}$& $\\cos{\\alpha}\/\\sin{\\beta}$ \\\\\n$\\xi^u_\\mathswitchr h$ & $\\cos{\\alpha}\/\\sin{\\beta}$ & $\\cos{\\alpha}\/\\sin{\\beta}$ & $\\cos{\\alpha}\/\\sin{\\beta}$ &$\\cos{\\alpha}\/\\sin{\\beta}$ \\\\\n$\\xi^d_\\mathswitchr h$& $\\cos{\\alpha}\/\\sin{\\beta}$ & $-\\sin{\\alpha}\/\\cos{\\beta}$ & $\\cos{\\alpha}\/\\sin{\\beta}$ &$-\\sin{\\alpha}\/\\cos{\\beta}$ \\\\\\hline\n$\\xi^l_\\mathswitchr {A_0}$ & $\\cot{\\beta}$ & $-\\tan{\\beta}$ &$-\\tan{\\beta}$ & $\\cot{\\beta}$ \\\\\n$\\xi^u_\\mathswitchr {A_0}$ & $\\cot{\\beta}$ & $\\cot{\\beta}$ & $\\cot{\\beta}$&$\\cot{\\beta}$ \\\\\n$\\xi^d_\\mathswitchr {A_0}$& $\\cot{\\beta}$ & $-\\tan{\\beta}$ & $\\cot{\\beta}$ & $-\\tan{\\beta}$ \\\\\\hline\n\\end{tabular} \n\\caption{The coupling strengths $\\xi^f$ of $\\mathswitchr H,\\mathswitchr h,\\mathswitchr {A_0}$ to the fermions relative to the SM value of $m_f\/v$, see Eq.~\\eqref{eq:ffH}. Note that the sign of $\\xi^f_\\mathswitchr {A_0}$ is defined relative to\nthe coupling of the Goldstone-boson field $G_0$ and the relation $\\beta_n=\\beta_c=\\beta$\nis used here.}\n\\label{tab:yukint}\n\\end{table}\nNote that we have used $\\beta_n=\\beta_c=\\beta$ in Eq.~\\refeq{eq:ffH} and \\refta{tab:yukint},\nwhich is most relevant in applications; the generalization to independent\n$\\beta_n$, $\\beta_c$, $\\beta$ is simple.\n\n\\section{The counterterm Lagrangian}\n\\label{sec:CTLagrangian}\n\nThe next step in calculating higher-order corrections is the renormalization of the theory. In this section we focus on electroweak corrections of $\\order{\\alpha_{\\mathrm{em}}}$. The QCD renormalization of the THDM is straightforward and completely analogous to the SM case, since all scalar degrees of freedom are colour singlets and do not interact strongly. \nIn the formulation of the basic Lagrangian in the previous section, we dealt with {\\it bare} \nparameters and fields. To distinguish those from {\\it renormalized} quantities, in the\nfollowing we indicate {\\it bare} quantities by subscripts ``0'' consistently.\nWe perform a multiplicative renormalization, i.e.\\ we split {\\it bare} quantities into\n{\\it renormalized} parts and corresponding counterterms,\nuse dimensional regularization,\nand allow for matrix-valued field renormalization constants in the case that there are several fields with the same quantum numbers.\nThe counterterm Lagrangian $\\delta \\mathcal{L}$ containing the full dependence on the renormalization constants can be split into several parts analogous to Eq.~\\eqref{eq:LTHDM},\n\\begin{align}\n\\label{eq:CTLagrangian}\n\\delta \\mathcal{L}_\\mathrm{THDM}\n=\\delta \\mathcal{L}_{\\mathrm{Gauge}}+\\delta \\mathcal{L}_{\\mathrm{Fermion}}+\\delta \\mathcal{L}_{\\mathrm{Higgs,kin}}-\\delta V_{\\mathrm{Higgs}}+\\delta \\mathcal{L}_{\\mathrm{Yukawa}},\n\\end{align}\nwhere the Higgs part of the Lagrangian $\\delta \\mathcal{L}_{\\mathrm{Higgs}}$ is split up into the kinetic part $\\delta \\mathcal{L}_{\\mathrm{Higgs,kin}}$ and the Higgs potential part $\\delta V_{\\mathrm{Higgs}}$.\nSince the gauge fixing is applied after renormalization, no gauge-fixing counterterms occur, and since ghost fields\noccur only in loop diagrams, a renormalization of the ghost sector is not necessary at NLO for the calculation of S-matrix elements. Though, for analyzing Slavnov--Taylor or Ward identities a complete renormalization procedure would be advisable. \n{Our renormalization procedure, thus, widely parallels the treatment described for the SM in \\cite{Denner:1991kt};\nan alternative variant that is based on the transformation of fields in the gauge eigenstate basis, as suggested\nfor the SM in \\citere{Bohm:1986rj} and for the MSSM in \\citere{Dabelstein:1994hb}, is described in\n\\citere{Altenkamp:Dis}.}\n\n\\subsection{Higgs potential}\n\\label{sec:CTpot}\nAccording to Eq.~\\eqref{eq:lambdapara}, the Higgs potential contains \n8 independent parameters which have to be renormalized. In addition, there are two vevs and two gauge couplings completing the set of input parameters of \\eqref{eq:basicpara}. \nWe have carried out different renormalization procedures and in this paper discuss the renormalization of the Lagrangian in the mass parameterization\n\\begin{enumerate}\n\\renewcommand{\\theenumi}{(\\alph{enumi}}\n\\renewcommand{\\labelenumi}{\\theenumi)}\n\\item with renormalization of the mixing angles,\n\\item alternatively,\ntaking the mixing angles as dependent parameters and applying the field transformation after renormalization.\n\\end{enumerate}\nAdditionally we have performed a renormalization of the basic parameters and a \nsubsequent transformation to renormalization constants and parameters of the mass parameter set similar to \nDabelstein in the MSSM~\\cite{Dabelstein:1994hb}%\n\\footnote{Details about this method can be found in \\citere{Altenkamp:Dis}.}. The latter has been used, after changing to our \nconventions for field and parameter renormalization, to check the counterterm Lagrangian. \nIn the following, we give a detailed description of method (a), while method~(b) is \nbriefly described in App.~\\ref{App:renormassa}.\n\n\\subsubsection{Renormalization of the mixing angles}\n\\label{sec:renmixangles}\n\n{In this section we show that the counterterms of mixing angles \nthat are not used to replace another free parameter of the theory \nare redundant in the sense that they can be absorbed by \nfield renormalization constants.}\nWe sketch the argument for generic scalar fields $\\varphi_{1},\\varphi_{2}$, which are transformed into fields $h_1, h_2$ \ncorresponding to mass eigenstates by a rotation by the angle~$\\theta$,\n\\begin{align}\n\\begin{pmatrix}\\varphi_{1,0} \\\\ \\varphi_{2,0}\\end{pmatrix}\n& = \\mathbf{R}_\\varphi(\\theta_{0}) \\begin{pmatrix}h_{1,0} \\\\ h_{2,0}\\end{pmatrix}\n=\\begin{pmatrix}c_{\\theta,0} & -s_{\\theta,0}\\\\ s_{\\theta,0}& c_{\\theta,0}\\end{pmatrix}\n\\begin{pmatrix}h_{1,0} \\\\ h_{2,0}\\end{pmatrix},\n\\label{eq:mixanglesrot}\n\\end{align}\nwhere we added subscripts ``0'' to indicate bare quantities.\nThe general argument can be applied to the neutral CP-even, the neutral \nCP-odd, and the charged Higgs fields of the THDM by replacing\n$h_1$, $h_2$ by $H$, $h$, or $G$, $A_0$ or $G^\\pm$, $H^\\pm$, respectively, and by substituting the angle $\\theta$ by \n$\\alpha$, $\\beta_n$, or $\\beta_c$.\nThe fields corresponding to mass eigenstates are renormalized using matrix-valued renormalization constants, so that the renormalization transformation reads\n\\begin{align}\n\\begin{pmatrix}h_{1,0} \\\\ h_{2,0} \\end{pmatrix}&=\\begin{pmatrix}1+\\frac{1}{2}\\delta Z_{\\mathswitchr h_1 \\mathswitchr h_1} &\\frac{1}{2} \\delta Z_{\\mathswitchr h_1 \\mathswitchr h_2} \\\\\\frac{1}{2} \\delta Z_{\\mathswitchr h_2 \\mathswitchr h_1} &1+\\frac{1}{2}\\delta Z_{\\mathswitchr h_2 \\mathswitchr h_2} \\end{pmatrix} \\begin{pmatrix}h_1 \\\\ h_2\\end{pmatrix},&&\n \\theta_{0} = \\theta + \\delta \\theta.\\label{eq:alpherendef}\n\\end{align}\nApplying this renormalization transformation to Eq.~\\eqref{eq:mixanglesrot} leads to\n\\begin{align}\n\\begin{pmatrix} \\varphi_{1,0} \\\\ \\varphi_{2,0}\\end{pmatrix}={}\n& \\bigg[ \\begin{pmatrix}c_{\\theta} & -s_{\\theta} \\\\ s_{\\theta}& c_{\\theta} \\end{pmatrix}\n\\begin{pmatrix}1+\\frac{1}{2}\\delta Z_{\\mathswitchr h_1\\mathswitchr h_1} & \\frac{1}{2}\\delta Z_{\\mathswitchr h_1\\mathswitchr h_2} \\\\ \\frac{1}{2}\\delta Z_{\\mathswitchr h_2\\mathswitchr h_1}& 1+\\frac{1}{2}\\delta Z_{\\mathswitchr h_2\\mathswitchr h_2} \\end{pmatrix}\n+\\begin{pmatrix}-s_{\\theta} & -c_{\\theta} \\\\ c_{\\theta}& -s_{\\theta} \\end{pmatrix} \\delta \\theta \\bigg]\n\\begin{pmatrix}h_1 \\\\ h_2\\end{pmatrix}\n\\nonumber\\\\\n{}={}& \\begin{pmatrix}c_{\\theta} & -s_{\\theta} \\\\ s_{\\theta}& c_{\\theta} \\end{pmatrix}\n\\begin{pmatrix}1+\\frac{1}{2}\\delta Z_{\\mathswitchr h_1\\mathswitchr h_1} &\\frac{1}{2}(\\delta Z_{\\mathswitchr h_1\\mathswitchr h_2}- 2 \\delta \\theta) \n\\\\ \n\t\\frac{1}{2}(\\delta Z_{\\mathswitchr h_2\\mathswitchr h_1}+ 2 \\delta \\theta)& 1+\\frac{1}{2}\\delta Z_{\\mathswitchr h_2\\mathswitchr h_2} \\end{pmatrix}\n\\begin{pmatrix}h_1 \\\\ h_2\\end{pmatrix}. \n\\label{eq:mixanglesrot2}\n\\end{align}\nOne can easily remove the dependence on the mixing angle \nwith a redefinition of the mixing \nfield renormalization constants by introducing\n\\begin{align}\n\\label{eq:redeffielren}\n \\delta \\tilde{Z}_{\\mathswitchr h_2\\mathswitchr h_1}= \\delta Z_{\\mathswitchr h_2\\mathswitchr h_1}+ 2 \\delta \\theta,\\qquad\n \\delta \\tilde{Z}_{\\mathswitchr h_1\\mathswitchr h_2}= \\delta Z_{\\mathswitchr h_1\\mathswitchr h_2}- 2 \\delta \\theta.\n \\end{align}\nThen, the Eq.~\\eqref{eq:mixanglesrot2} reads\n\\begin{align}\n\\begin{pmatrix} \\varphi_{1,0} \\\\ \\varphi_{2,0}\\end{pmatrix}=\n& \\begin{pmatrix}c_{\\theta} & -s_{\\theta} \\\\ s_{\\theta}& c_{\\theta} \\end{pmatrix}\n\\begin{pmatrix}1+\\frac{1}{2}\\delta Z_{\\mathswitchr h_1\\mathswitchr h_1} &\\frac{1}{2}\\delta \\tilde{Z}_{\\mathswitchr h_1\\mathswitchr h_2} \\\\ \n\t\\frac{1}{2}\\delta \\tilde{Z}_{\\mathswitchr h_2\\mathswitchr h_1}& 1+\\frac{1}{2}\\delta Z_{\\mathswitchr h_2\\mathswitchr h_2} \\end{pmatrix}\n\\begin{pmatrix}h_1 \\\\ h_2\\end{pmatrix}.\n\\end{align}\nObviously,\nthe dependence on $\\delta \\theta$ can always be removed from the Lagrangian \nby a redefinition of the field renormalization constants. \nAs a simple shift of the mixing field renormalization constant is performing the task, the renormalization of the mixing angle $\\theta$ can be seen as an additional field renormalization \n(as it is done, e.g., in Ref.~\\cite{Aoki2009}). \nThis argument is general and holds for any renormalization condition on $\\theta$.\nWithout loss of generality one can even assume that such a redefinition has already been performed and set $\\delta \\theta=0$ from the beginning, as done in method (b) in App.~\\ref{App:renormassa}.\nOf course, the bookkeeping of counterterms depends on the way $\\delta \\theta$ is treated.\nThis can be seen by considering the mass term of the potential. \nThe general bare mass term can be written using the rotation matrix $\\mathbf{R_\\varphi}(\\theta_{0})$ as\n\\begin{align}\nV_{\\mathswitchr h_1\\mathswitchr h_2} {}={} \n&\\frac{1}{2} \\,(h_{1,0}, h_{2,0})\\, \\mathbf{R}_\\varphi^{\\mathrm{T}}(\\theta_{0}) \\mathbf{M}_{\\varphi,0} \n\\mathbf{R}_\\varphi(\\theta_{0}) \\begin{pmatrix} h_{1,0}\\\\ h_{2,0} \\end{pmatrix}\n\\nonumber\\\\\n{}={} &\\frac{1}{2} \\, (h_{1,0}, h_{2,0})\\, \n\\mathbf{R}_\\varphi^{\\mathrm{T}}(\\theta+\\delta \\theta) \n\\left(\\mathbf{M_\\varphi}+\\delta \\mathbf{M_\\varphi}\\right) \n\\mathbf{R_\\varphi}(\\theta+\\delta \\theta) \n\\begin{pmatrix} h_{1,0}\\\\ h_{2,0} \\end{pmatrix}.\n\\end{align}\nThis expression can be expanded in terms of renormalized and counterterm contributions, yielding\n\\begin{align}\nV_{\\mathswitchr h_1\\mathswitchr h_2}=\n\\frac{1}{2} \\, (h_{1,0}, h_{2,0})\\, \\begin{pmatrix} \n\\hspace{-0pt}M_{\\mathswitchr h_1}^2+\\delta M_{\\mathswitchr h_1}^2 & \n\\hspace{-5pt}\\delta \\theta (M_{\\mathswitchr h_2}^2-M_{\\mathswitchr h_1}^2)+f_{\\theta}(\\{ \\delta p\\}) \\\\ \\delta \\theta (M_{\\mathswitchr h_2}^2-M_{\\mathswitchr h_1}^2)+f_{\\theta}(\\{ \\delta p\\}) & \n\\hspace{-0pt}M_{\\mathswitchr h_2}^2+\\delta M_{\\mathswitchr h_2}^2\\end{pmatrix} \\begin{pmatrix} h_{1,0}\\\\ h_{2,0} \\end{pmatrix},\n\\end{align}\nwhere we obtain off-diagonal terms from the counterterm $\\delta\\theta$\nof the mixing angle and from the renormalization of the mass matrix. The latter contribution depends on the independent counterterms $\\{ \\delta p\\}$ and is abbreviated by $f_{\\theta}(\\{ \\delta p\\})$.\nAt NLO, the mixing entry of the mass matrix reads\n\\begin{align}\n\\delta M_{\\mathswitchr h_1\\mathswitchr h_2}^2= \\delta \\theta (M_{\\mathswitchr h_2}^2-M_{\\mathswitchr h_1})+f_{\\theta}(\\{ \\delta p\\}), \\label{eq:NLOmassmixHh}\n\\end{align}\n{and since the renormalization of the {\\it redundant} mixing angle can be chosen freely,} \ncounterterm contributions in the Lagrangian can be shifted arbitrarily from mixing terms to mixing-angle counterterms.\n{Note that $f_{\\theta}(\\{ \\delta p\\})$ does not change\nby such redistributions, since it is fixed by the remaining renormalization \nconstants.}\n\n\\subsubsection{Renormalization with a diagonal mass matrix -- version (a)}\n\\label{sec:renormassb}\n\nIn this prescription, we use the angle $\\alpha$ as an independent parameter instead of $\\lambda_3$. \n{To define $\\alpha$ at NLO, we demand that \nthe mass matrix of the CP-even Higgs bosons (in the Lagrangian), written in terms of bare fields,\nis diagonal at all orders (i.e.\\ in terms of bare or renormalized parameters).}\nEquation~\\eqref{eq:Lambda3NLO} with\n\\begin{align}\\label{eq:MHhchoice2}\nM_{\\mathrm{Hh},0}^2 = 0 + \\delta M_{\\mathrm{Hh}}^2=0, \\qquad M_{\\mathrm{Hh}}^2 = 0\n\\end{align}\nthen defines the parameter $\\alpha$. \n{Note that (momentum-dependent)\nloop diagrams tend to destroy the diagonality \nof the matrix-valued two-point functions (inverse propagators)\nin the effective action as well.\nBelow,\nthe field renormalization will be chosen to compensate those loop effects \nat the mass shells of the propagating particles.}\n{It is relation \\refeq{eq:MHhchoice2} that distinguishes \n$\\alpha$ from the case of a redundant mixing angle, such as $\\theta$ in the\nprevious section, which can be chosen freely or absorbed by field\nrenormalization constants.}\nThe relation between $\\alpha$ and $\\lambda_3$ of Eq.~\\eqref{eq:Lambda3NLO} is the same for bare and renormalized quantities \nand can be used to eliminate $\\lambda_3$ from the theory.\nFor each of the independent parameters of the mass parameter set \\eqref{eq:physparaNLO} we apply the renormalization transformation\n\\begin{align}\nM_{\\mathswitchr H,0}^2&=\\mathswitch {M_\\PH}^2+\\delta \\mathswitch {M_\\PH}^2,&M_{\\mathswitchr h,0}^2&= \\mathswitch {M_\\Ph}^2+\\delta \\mathswitch {M_\\Ph}^2,&M_{\\mathswitchr {A_0},0}^2 &= \\mathswitch {M_\\PAO}^2+\\delta \\mathswitch {M_\\PAO}^2,\\nonumber\\\\\nM_{\\mathswitchr {H^+},0}^2&= \\mathswitch {M_\\PHP}^2+\\delta \\mathswitch {M_\\PHP}^2,&M_{\\mathswitchr W,0}^2&= \\mathswitch {M_\\PW}^2+\\delta \\mathswitch {M_\\PW}^2,&M_{\\mathswitchr Z,0}^2&= \\mathswitch {M_\\PZ}^2+\\delta \\mathswitch {M_\\PZ}^2, \\nonumber\\\\\ne_0&= e+\\delta e,&\\lambda_{5,0} &= \\lambda_5+\\delta \\lambda_5&\\alpha_0 &= \\alpha+\\delta \\alpha,\\nonumber\\\\\n\\beta_0&= \\beta+\\delta \\beta,&t_{\\mathswitchr H,0}&= 0+\\delta t_\\mathswitchr H,&t_{\\mathswitchr h,0}&= 0+\\delta t_\\mathswitchr h, \\label{eq:physpararen2}\n\\end{align}\nso that the 12 parameter renormalization constants are\n\\begin{align}\n\\{\\delta p_\\mathrm{mass}\\}=\\{\\delta \\mathswitch {M_\\PH}^2, \\delta \\mathswitch {M_\\Ph}^2, \\delta \\mathswitch {M_\\PAO}^2, \\delta \\mathswitch {M_\\PHP}^2, \\delta \\mathswitch {M_\\PW}^2, \\delta \\mathswitch {M_\\PZ}^2, \\delta e, \\delta \\lambda_5, \\delta \\alpha,\\delta \\beta,\\delta t_\\mathswitchr H, \\delta t_\\mathswitchr h\\},\n\\label{eq:renconstsR}\n\\end{align}\ncorresponding to $\\{p_\\mathrm{mass}\\}$. \nIn this part we describe the commonly used tadpole renormalization (which is gauge dependent) and describe \na gauge-independent scheme in Sect.~\\ref{sec:FJscheme}.\n\nThe higher-order corrections of the mixing angles $\\beta_n$ and $\\beta_c$ are irrelevant according to Sect.~\\ref{sec:renmixangles} and we can choose\n\\begin{align}\n\\beta_{n,0}=\\beta_{c,0}= \\beta_0 = \\beta+\\delta \\beta,\n\\label{eq:mixanglect}\n\\end{align}\nwhich defines the mixing terms uniquely and ensures that the angles $\\beta_n,\\beta_c$, and $\\beta$ do not have to be distinguished at any order.\nFrom these conditions, we can compute the mass mixing terms from Eq.~\\eqref{eq:goldstonemassmixing} and Eq.~\\eqref{eq:chargedgoldstonemassmixing} to\n\\begin{align}\n M_{\\mathswitchr {G_0}\\mathswitchr {A_0},0}^2 &= 0+\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2\n=-e\\frac{ \n\\delta t_\\mathswitchr H s_{\\alpha -\\beta }+\n\\delta t_\\mathswitchr h c_{\\alpha -\\beta }\n}{2 \\mathswitch {M_\\PW} \\mathswitch {s_{\\scrs\\PW}}},&\\nonumber\\\\\n M_{\\mathswitchr {G}\\mathswitchr {H^+},0}^2 &= 0+\\delta M_{\\mathswitchr {G}\\mathswitchr {H^+}}^2\n=-e\\frac{ \n\\delta t_\\mathswitchr H s_{\\alpha -\\beta }+\n\\delta t_\\mathswitchr h c_{\\alpha -\\beta }\n}{2 \\mathswitch {M_\\PW} \\mathswitch {s_{\\scrs\\PW}}}. \\label{eq:Mabchoice2}\n\\end{align}\nThe field re\\-nor\\-ma\\-li\\-zation is performed \nfor each field corresponding to mass eigenstates,\n\\begin{align}\n\\begin{pmatrix}H_0 \\\\ h_0 \\end{pmatrix}&=\\begin{pmatrix}1+\\frac{1}{2}\\delta Z_\\mathswitchr H &\\frac{1}{2} \\delta Z_{\\mathswitchr H\\mathswitchr h} \\\\\\frac{1}{2} \\delta Z_{\\mathswitchr h\\mathswitchr H} &1+\\frac{1}{2}\\delta Z_{\\mathswitchr h} \\end{pmatrix} \\begin{pmatrix}H \\\\ h\\end{pmatrix},\\nonumber\\\\\n\\begin{pmatrix}G_{0,0} \\\\ A_{0,0} \\end{pmatrix}&=\\begin{pmatrix}1+\\frac{1}{2}\\delta Z_{\\mathswitchr {G_0}} &\\frac{1}{2}\\delta Z_{\\mathswitchr {G_0}\\mathswitchr {A_0}} \\\\\\frac{1}{2}\\delta Z_{\\mathswitchr {A_0}\\mathswitchr {G_0}} &1+\\frac{1}{2}\\delta Z_\\mathswitchr {A_0} \\end{pmatrix} \\begin{pmatrix}G_0 \\\\ A_0\\end{pmatrix},\\nonumber\\\\\n\\begin{pmatrix}G^\\pm_0 \\\\ H^\\pm_0 \\end{pmatrix}&=\\begin{pmatrix}1+\\frac{1}{2}\\delta Z_{\\mathrm{G}^+} &\\frac{1}{2}\\delta Z_{\\mathrm{GH}^+} \\\\\\frac{1}{2}\\delta Z_{\\mathrm{HG}^+} &1+\\frac{1}{2}\\delta Z_{\\mathswitchr H^+} \\end{pmatrix} \\begin{pmatrix}G^\\pm \\\\ H^\\pm\\end{pmatrix},\\label{eq:physfieldren}\n\\end{align}\nwith the field renormalization constants $\\delta Z_\\mathswitchr H$, $\\delta Z_{\\mathswitchr H\\mathswitchr h}$, $\\delta Z_{\\mathswitchr h\\mathswitchr H}$, $\\delta Z_{\\mathswitchr h}$, $\\delta Z_{\\mathswitchr {G_0}}$, \n$\\delta Z_{\\mathswitchr {G_0}\\mathswitchr {A_0}}$, $\\delta Z_{\\mathswitchr {A_0}\\mathswitchr {G_0}}$, $\\delta Z_\\mathswitchr {A_0}$, $\\delta Z_{\\mathrm{G}^+}$, $\\delta Z_{\\mathrm{GH}^+}$, $\\delta Z_{\\mathrm{HG}^+}$, and $\\delta Z_{\\mathswitchr H^+}$ for the CP-even, the CP-odd, and the charged Higgs fields.\nWe denote the complete set of parameter and field renormalization constants with $\\{\\delta \\mathbf{R}_\\mathrm{mass}\\}$. All renormalization constants \nare of $\\order{\\alpha_{\\mathrm{em}}}$, i.e.\\ all contributions of $\\order{\\alpha_{\\mathrm{em}}^2}$ or higher are omitted. \n\n\\begin{sloppypar}\nApplying the renormalization transformations~\\refeq{eq:physpararen2}, \\refeq{eq:physfieldren} \nto the bare potential of Eq.~\\eqref{eq:NLObarepot} and linearizing in the renormalization constants results in \n$V(\\{p_\\mathrm{mass}\\})+\\delta V(\\{p_\\mathrm{mass}\\},\\{\\delta \\mathbf{R}_\\mathrm{mass}\\})$\nwith the LO potential as in Eq.~\\eqref{eq:NLObarepot}, but with renormalized quantities and the counterterm potential of $\\order{\\alpha_\\mathrm{em}}$,\n\\begin{align}\n\\delta V(\\{p_\\mathrm{mass}\\},\\{\\delta \\mathbf{R}_\\mathrm{mass}\\})=&- \\delta t_\\mathswitchr H H -\\delta t_\\mathswitchr h h \\nonumber\\\\\n&+\\frac{1}{2} \\left(\\delta \\mathswitch {M_\\PH}^2 +\\delta Z_{\\mathswitchr H} \\mathswitch {M_\\PH}^2\\right) H^2 +\\frac{1}{2} \\left(\\delta \\mathswitch {M_\\Ph}^2 +\\delta Z_{\\mathswitchr h} \\mathswitch {M_\\Ph}^2\\right) h^2 \\nonumber\\\\\n&+\\frac{1}{2} \\left(\\delta \\mathswitch {M_\\PAO}^2 +\\delta Z_{\\mathswitchr {A_0}} \\mathswitch {M_\\PAO}^2\\right) A_0^2 + \\left(\\delta \\mathswitch {M_\\PHP}^2+\\delta Z_{\\mathswitchr {H^+}} \\mathswitch {M_\\PHP}^2\\right)H^+ H^-\\nonumber\\\\\n&+ \\frac{e}{4 \\mathswitch {M_\\PW} \\mathswitch {s_{\\scrs\\PW}}} \\left(\n-\\delta t_\\mathswitchr H c_{\\alpha-\\beta}\n+\\delta t_\\mathswitchr h s_{\\alpha-\\beta}\n\\right) (G_0^2 + 2 G^+ G^-)\\nonumber\\\\\n&+ \\frac{1}{2}\\left( \\mathswitch {M_\\PH}^2\\delta Z_{\\mathswitchr H\\mathswitchr h}+\\mathswitch {M_\\Ph}^2 \\delta Z_{\\mathswitchr h\\mathswitchr H}\\right) H h\\nonumber\\\\\n&+\\frac{1}{2} \\left( \\mathswitch {M_\\PAO}^2 \\delta Z_{\\mathswitchr {A_0}\\mathswitchr {G_0}} \n+ 2\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2 \\right) G_0 A_0\n\\nonumber\\\\\n&+ \\frac{1}{2} \\left( \\mathswitch {M_\\PHP}^2 \\delta Z_{\\mathswitchr H\\mathswitchr {G^+}} \n+2\\delta M_{\\mathswitchr {G}\\mathswitchr H+}^2 \\right)(H^+ G^-+G^+ H^-)\n\\nonumber \\\\\n&+ \\textrm{interaction terms}.\\label{eq:dVH2p}\n\\end{align}\nThe interaction terms are derived in\nthe same way, but they are very lengthy and not shown here. \n\\end{sloppypar}\n\nThe prescription for the field renormalization~\\eqref{eq:physfieldren} is non-minimal in the sense that a renormalization of the doublets with two renormalization constants,\n\\begin{align}\n\\Phi_{1,0} = Z_{\\mathswitchr H_1}^{1\/2} \\Phi_1& \\textstyle = \\Phi_1 \\left(1+\\frac{1}{2} \\delta Z_{\\mathswitchr H_1}\\right),\\nonumber\\\\\n\\Phi_{2,0} = Z_{\\mathswitchr H_2}^{1\/2} \\Phi_2& \\textstyle = \\Phi_2 \\left(1+\\frac{1}{2} \\delta Z_{\\mathswitchr H_2}\\right),\\label{eq:inputfieldren}\n\\end{align}\nactually would be\nsufficient to cancel the UV divergences. However, the prescription with matrix-valued renormalization constants allows us to renormalize each field on-shell. The UV-divergent parts of \nthe renormalization constants in Eq.~\\refeq{eq:physfieldren}\ncannot be independent, and relations between the UV-divergent parts of the two prescriptions exist. They can be obtained by applying the renormalization prescription \\eqref{eq:inputfieldren} to the left-hand side and \\eqref{eq:physfieldren} to the right-hand side of Eqs.~\\eqref{eq:rotations}, transforming thereafter the interaction states on the left-hand side to mass eigenstates and comparing both sides.\nThis results in\n\\begin{align}\n\\delta Z_\\mathswitchr h\\big|_\\mathrm{UV}&=s_{\\alpha}^2\\;\\delta Z_{\\mathrm{H}_1}\\big|_\\mathrm{UV}+c_{\\alpha}^2\\; \\delta Z_{\\mathrm{H}_2}\\big|_\\mathrm{UV},\\nonumber\\\\\n\\delta Z_\\mathswitchr H\\big|_\\mathrm{UV}&=c_{\\alpha}^2\\; \\delta Z_{\\mathrm{H}_1}\\big|_\\mathrm{UV}+s_{\\alpha}^2\\; \\delta Z_{\\mathrm{H}_2}\\big|_\\mathrm{UV},\\nonumber\\\\\n\\delta Z_{\\mathswitchr H\\mathswitchr h}\\big|_\\mathrm{UV}&=s_{\\alpha} c_{\\alpha}\\; \\left(-\\delta Z_{\\mathrm{H}_1}\\big|_\\mathrm{UV}+\\delta Z_{\\mathrm{H}_2}\\big|_\\mathrm{UV}\\right)+2 \\delta \\alpha\\big|_\\mathrm{UV},\\nonumber\\\\\n\\delta Z_{\\mathswitchr h\\mathswitchr H}\\big|_\\mathrm{UV}&=s_{\\alpha} c_{\\alpha}\\; \\left(-\\delta Z_{\\mathrm{H}_1}\\big|_\\mathrm{UV}+\\delta Z_{\\mathrm{H}_2}\\big|_\\mathrm{UV}\\right)-2 \\delta \\alpha\\big|_\\mathrm{UV},\\nonumber\\\\\n\\delta Z_\\mathswitchr {A_0}\\big|_\\mathrm{UV}&=\\delta Z_{\\mathswitchr H^+}\\big|_\\mathrm{UV}=s_{\\beta}^2\\;\\delta Z_{\\mathrm{H}_1}\\big|_\\mathrm{UV}+c_{\\beta}^2\\; \\delta Z_{\\mathrm{H}_2}\\big|_\\mathrm{UV},\\nonumber\\\\\n\\delta Z_\\mathswitchr {G_0}\\big|_\\mathrm{UV}&=\\delta Z_{G^+}\\big|_\\mathrm{UV}=c_{\\beta}^2\\; \\delta Z_{\\mathrm{H}_1}\\big|_\\mathrm{UV}+s_{\\beta}^2\\; \\delta Z_{\\mathrm{H}_2}\\big|_\\mathrm{UV},\\nonumber\\\\\n\\delta Z_{\\mathswitchr {G_0}\\mathswitchr {A_0}}\\big|_\\mathrm{UV}&=\\delta Z_{\\mathswitchr {G}\\mathswitchr {H^+}}\\big|_\\mathrm{UV}=s_{\\beta} c_{\\beta}\\; \\left(-\\delta Z_{\\mathrm{H}_1}\\big|_\\mathrm{UV}+\\delta Z_{\\mathrm{H}_2}\\big|_\\mathrm{UV}\\right)+2 \\delta \\beta\\big|_\\mathrm{UV},\\nonumber\\\\\n\\delta Z_{\\mathswitchr {A_0}\\mathswitchr {G_0}}\\big|_\\mathrm{UV}&=\\delta Z_{\\mathswitchr H\\mathswitchr {G^+}}\\big|_\\mathrm{UV}=s_{\\beta} c_{\\beta}\\; \\left(-\\delta Z_{\\mathrm{H}_1}\\big|_\\mathrm{UV}+\\delta Z_{\\mathrm{H}_2}\\big|_\\mathrm{UV}\\right)-2 \\delta \\beta\\big|_\\mathrm{UV}.\\label{eq:fieldrenconstrel}\n\\end{align}\nWe will use these relations to derive UV-divergent parts for specific renormalization constants in Sect.~\\ref{sec:diffrenschemes}. In App.~\\ref{App:renormassa} we discuss a different choice of the mixing angles which is suited for \nthe renormalization with $\\lambda_3$ as an independent parameter. \n\n\\subsection{The Higgs kinetic part}\n\\label{sec:Hkinren}\n\nAfter expressing the Higgs kinetic term\n\\begin{align}\n \\mathcal{L}_\\mathrm{H,kin}=(D_\\mu\\Phi_1)^\\dagger (D^\\mu \\Phi_1)+(D_\\mu\\Phi_2)^\\dagger (D^\\mu \\Phi_2)\n\\end{align}\nin terms of\nbare physical fields, mixing angles, and parameters, one can apply the renorma\\-lization transformations \\eqref{eq:physpararen2} and~\\eqref{eq:physfieldren} to obtain the counterterm part of the kinetic Lagrangian which introduces \nscalar--vector mixing terms. The explicit terms are stated in App.~\\ref{App:SV-terms}.\n\n\\subsection{Fermionic and gauge parts}\n\\label{sec:fermgauge}\nSince the THDM extension of the SM does not affect the gauge and the fermion parts of the Lagrangian, \nthe renormalization of these parts is identical to the SM case. \n{It is described in detail in Ref.~\\cite{Denner:1991kt} in \nBHS convention, which is included in the standard implementation of the \\textsc{FeynArts} package~\\cite{Hahn2001}.\nOther renormalization prescriptions\ncan, e.g., be found in \\citeres{Denner:2004bm,Kniehl:2009kk}.}\nTherefore, we here do not repeat the renormalization procedure of the CKM matrix, which does not change in the \ntransition from the SM to the THDM, and spell out the renormalization of the fermionic parts only for the\ncase where the CKM matrix is set to the unit matrix, i.e.\\ $V_{ij}=\\delta_{ij}$.\nThe transformation of the left- and right-handed fermions and of the gauge-boson fields are\n\\begin{align}\n f^\\sigma_{i,0} & = \\left(1 + \\textstyle\\frac{1}{2}\\delta Z^{f,\\sigma}_i\\right) f^\\sigma_i,\\qquad \n\\rlap{$f=\\nu,l,u,d, \\quad \\sigma={\\mathrm{L}},{\\mathrm{R}}, \\quad i=1,2,3,$}\\\\\n \\begin{pmatrix}Z_0 \\\\ A_\\mathrm{bare}\\end{pmatrix}&=\\begin{pmatrix}1+\\frac{1}{2}\\delta Z_{\\mathswitchr Z\\PZ} &\\frac{1}{2} Z_{\\mathswitchr Z{\\mathrm{PA}}} \\\\\\frac{1}{2} \\delta Z_{{\\mathrm{PA}}\\mathswitchr Z} &1+\\frac{1}{2}\\delta Z_{{\\mathrm{PA}}\\PA} \\end{pmatrix} \\begin{pmatrix}Z \\\\ A\\end{pmatrix}, &\nW^\\pm_0 &= \\left(1 + \\textstyle\\frac{1}{2} \\delta Z_\\mathswitchr W\\right) W^\\pm,\n\\end{align}\nwhere the bare photon field is denoted $A_\\mathrm{bare}$ to distinguish it from the \nneutral CP-odd field~$A_0$.\nMixing between left-handed up- and down-type fermions does not occur, owing to charge conservation. Inserting this into the Lagrangian directly delivers the renormalized and the counterterm Lagrangians.\n\n\\subsection{Yukawa part}\n\\label{sec:YukRen}\n\nThe renormalization of the Yukawa sector is straightforward in Type I, II, lepton-specific, and flipped models and can be done by taking the Lagrangian of Eq.~\\eqref{eq:ffH}, \nreplacing the vev $v$, and applying the renormalization transformations of \nSect.~\\ref{sec:renormassb}, as well as a renormalization of the fermion masses,\n\\begin{align}\n m_{f,i,0} &= m_{f,i} + \\delta m_{f,i}.\n \\end{align}\nThe corresponding counterterm couplings are stated in App.~\\ref{App:Yuk-CT}.\n\n\\section{Renormalization conditions}\n\\label{sec:renconditions}\n\nThe renormalization constants are fixed using on-shell conditions \nfor all parameters that are accessible by experiments. \nHowever, not all parameters of the THDM correspond to measurable quantities, so that we renormalize three parameters of the Higgs sector in the $\\overline{\\mathrm{MS}}$ scheme, where the renormalization constants only contain the standard UV divergence\n\\begin{align}\n \\Delta_\\mathrm{UV}=\\frac{2}{4-D}-\\gamma_\\mathrm{E}+\\ln{4\\pi}=\\frac{1}{\\epsilon}-\\gamma_\\mathrm{E}+\\ln{4\\pi}\n\\end{align}\nin $D=4-2 \\epsilon$ dimensions and with the Euler--Mascheroni constant $\\gamma_\\mathrm{E}$.\nIn Sect.~\\ref{sec:diffrenschemes}, four different options, resulting in four different renormalization schemes, are presented.\nAn overview over the renormalization constants introduced in the previous section is shown in Tab.~\\ref{tab:renconsts}. In the following, we adapt the notation of Ref.~\\cite{Denner:1991kt}, i.e.\\ we use the same symbols for the renormalized and the corresponding unrenormalized \nGreen function, self-energies, etc., but denoting the renormalized quantities with a caret.\n\n\\begin{table}\n \\centering\n\\begin{tabular}{|lrl|}\\hline\nParameters:&&\\\\\\hline\n&EW (3): & $\\delta \\mathswitch {M_\\PZ}^2,\\delta \\mathswitch {M_\\PW}^2,\\delta e, (\\delta \\mathswitch {c_{\\scrs\\PW}},\\delta \\mathswitch {s_{\\scrs\\PW}})$\\\\\n&fermion masses (9): & $\\delta m_{f,i},\\qquad f=l,u,d,\\quad i=1,2,3$\\\\%\\delta m_{u,i},\\delta m_{d,i}$\\\\\n&Higgs masses (4): & $\\delta \\mathswitch {M_\\PH}^2,\\delta \\mathswitch {M_\\Ph}^2,\\delta \\mathswitch {M_\\PAO}^2,\\delta \\mathswitch {M_\\PHP}^2$\\\\\n&Higgs potential (3): & $\\delta \\lambda_3 \\text{ or } \\delta \\alpha,\\delta \\lambda_5, \\delta \\beta$\\\\\n&tadpoles (2): & $\\delta t_\\mathswitchr H,\\delta t_\\mathswitchr h$\\\\\\hline\nFields:&&\\\\\\hline\n&EW (5): & $\\delta Z_\\mathswitchr W,\\delta Z_{\\mathswitchr Z\\PZ},\\delta Z_{\\mathswitchr Z{\\mathrm{PA}}},\\delta Z_{{\\mathrm{PA}}\\mathswitchr Z},\\delta Z_{{\\mathrm{PA}}\\PA}$\\\\\n&left-handed fermions (12): & $\\delta Z^{f,{\\mathrm{L}}}_i,\\qquad f=\\nu,l,u,d,\\quad i=1,2,3$ \\\\\n&right-handed fermions (9): & $\\delta Z^{f,{\\mathrm{R}}}_i,\\qquad f=l,u,d,\\quad i=1,2,3$ \\\\\n&Higgs (12): & $\\delta Z_{\\mathswitchr H},\\delta Z_{\\mathswitchr H\\mathswitchr h},\\delta Z_{\\mathswitchr h\\mathswitchr H},\\delta Z_\\mathswitchr h$\\\\\n&& $\\delta Z_{\\mathswitchr {A_0}},\\delta Z_{\\mathswitchr {A_0}\\mathswitchr {G_0}},\\delta Z_{\\mathswitchr {G_0}\\mathswitchr {A_0}},\\delta Z_\\mathswitchr {G_0}$\\\\\n&&$\\delta Z_{\\mathswitchr {H^+}},\\delta Z_{\\mathswitchr H\\mathswitchr {G^+}},\\delta Z_{\\mathswitchr {G}\\mathswitchr {H^+}},\\delta Z_{\\mathswitchr {G^+}}$\\\\\\hline\n\\end{tabular} \n \\caption{The renormalization constants used to describe the THDM, separated into sectors of parameter and field renormalization. The renormalization constants in parentheses are not independent, but useful for a better bookkeeping. The numbers in parentheses are the numbers of independent renormalization constants. In total there are 38 field and 20 parameter renormalization constants to fix. }\n\\label{tab:renconsts}\n\\end{table}%\n\n\\subsection{On-shell renormalization conditions}\n\\label{sec:onshellrenconditions}\n\n\\subsubsection{Higgs sector}\n\n\\paragraph{Tadpoles:}\nWe start with the \n(irreducible) renormalized one-point vertex function\n\\begin{align}\n\\hat\\Gamma^{\\mathswitchr H,\\mathswitchr h}=\n{\\mathrm{i}}\\hat{T}^{\\mathswitchr H,\\mathswitchr h}=\\,\\parbox{2.5cm}{\\includegraphics[scale=1]{.\/diagrams\/tadpole1.eps}}\n\\hspace{20pt}.\n\\end{align}\nAt NLO the renormalized tadpole $\\hat{T}$ consists of a counterterm contribution $\\delta t$ and an unrenormalized one-loop irreducible one-point vertex function $T$\nresulting from the diagrams shown in Fig.~\\ref{fig:tadpoles}. \n\\begin{figure}\n\\vspace{-30pt}\n\\input{diagrams\/tadpoles}\n\\vspace{-40pt}\n\\caption{Generic tadpole diagrams. There is one diagram for each massive fermion $(f)$, \nscalar $(S)$, gauge-boson $(V)$, and ghost field $(u)$.}\n\\label{fig:tadpoles}\n\\end{figure}\nIn the conventional, but gauge-dependent tadpole treatment one demands that these two contributions cancel each other,\n\\begin{align}\n\\hat{T}_\\mathswitchr H&=\\delta t_\\mathswitchr H+T_\\mathswitchr H=0,&\\hat{T}_\\mathswitchr h&=\\delta t_\\mathswitchr h+T_\\mathswitchr h=0,\n\\label{eq:tadpoleren}\n\\end{align}\nwhich means that explicit tadpole diagrams can be omitted from the \nset of one-loop diagrams for any process. \nHowever, as a remnant of the tadpole diagrams the tadpole counter\\-terms \nappear also in various coupling \ncounterterms and need to be calculated. \nIt should be noted \nthat the condition on the tadpoles does not affect physical observables \nas long as physically equivalent renormalization conditions are imposed on\nthe input parameters. This is, in particular, the case for on-shell renormalization,\nwhere input parameters are tied to measurable quantities.\nThat means, changing the tadpole renormalization condition shifts contributions between \nGreen functions and counterterms and merely changes the bookkeeping,\nbut the dependence of predicted observables on renormalized input parameters remains the same.\nThe situation changes if an {$\\overline{\\mathrm{MS}}$}{} renormalization condition is used, where the counterterm\nis not fixed by a measurable quantity, \nbut by a divergence in a specific Green function, so that\nthe gauge-dependent tadpole terms can affect the relation between renormalized input parameters and\nobservables.\nThe gauge-independent treatment of tadpole contributions is based on a different renormalization condition and discussed in Sect.~\\ref{sec:FJscheme}. \n\\vspace{10pt}\n\n\\paragraph{Scalar self-energies:}\nFor scalars, the irreducible two-point functions with momentum transfer $k$ are\n\\begin{align}\n\\hat{\\Gamma}^{{ab}}(k)&=\\quad \\parbox{3.0cm}{\n\\includegraphics[scale=1]{.\/diagrams\/SE1.eps}} \\quad\n=\\mathrm{i}\\delta_{{ab}} (k^2-M_{a}^2)+\\mathrm{i} \\hat{\\Sigma}^{{ab}}(k),\n\\label{eq:two-def}\n\\end{align}\nwhere both fields $a,b$ are incoming and $a,b=H,h,A_0,G_0,H^\\pm,G^\\pm$. The first term is the LO two-point vertex function, while the functions $\\hat{\\Sigma}^{{ab}}$ are the renormalized self-energies containing loop diagrams and counterterms. Generic diagrams contributing to the self-energies are shown in Fig.~\\ref{fig:selfenergies}. \n\\begin{figure}\n\\vspace{-10pt}\n\\input{diagrams\/selfenergies}\n\\vspace{-20pt}\n\\caption{Generic self-energy diagrams for the heavy, neutral CP-even Higgs self-energy, for other scalar self-energies the diagrams are analogous. Only massive particles contribute.}\n\\label{fig:selfenergies}\n\\end{figure}\nMixing occurs only between $H$ and $h$,\nbetween $A_0$ and $G_0$, and between $H^\\pm$ and $G^\\pm$.\nFor the neutral CP-even fields we obtain\n\\begin{subequations}\n\\label{eq:neutralSE}\n\\begin{align}\n\\hat{\\Sigma}^{\\mathswitchr h\\Ph}(k^2)&= \\Sigma^{\\mathswitchr h\\Ph}(k^2) + \\delta Z_\\mathswitchr h(k^2-\\mathswitch {M_\\Ph}^2) - \\delta \\mathswitch {M_\\Ph}^2,\\\\\n\\hat{\\Sigma}^{\\mathswitchr H\\PH}(k^2)&= \\Sigma^{\\mathswitchr H\\PH}(k^2) + \\delta Z_\\mathswitchr H(k^2-\\mathswitch {M_\\PH}^2) - \\delta \\mathswitch {M_\\PH}^2,\\\\\n\\hat{\\Sigma}^{\\mathswitchr H\\mathswitchr h}(k^2)&= \\Sigma^{\\mathswitchr H\\mathswitchr h}(k^2) + \\frac{1}{2}\\delta Z_{\\mathswitchr H\\mathswitchr h}(k^2-\\mathswitch {M_\\PH}^2) +\\frac{1}{2}\\delta Z_{\\mathswitchr h\\mathswitchr H}(k^2-\\mathswitch {M_\\Ph}^2) - \\delta M_{\\mathswitchr H\\mathswitchr h}^2, \\label{eq:twopointHh}\n\\end{align}\n\\end{subequations}\nand for the CP-odd fields\n\\begin{subequations}\n\\label{eq:CPoddSE}\n\\begin{align}\n\\hat{\\Sigma}^{\\mathswitchr {A_0}\\PAO}(k^2)&= \\Sigma^{\\mathswitchr {A_0}\\PAO}(k^2) + \\delta Z_{\\mathswitchr {A_0}}(k^2-\\mathswitch {M_\\PAO}^2) - \\delta \\mathswitch {M_\\PAO}^2,\\\\\n\\hat{\\Sigma}^{\\mathswitchr {G_0}\\PGO}(k^2)&= \\Sigma^{\\mathswitchr {G_0}\\PGO}(k^2) + \\delta Z_\\mathswitchr {G_0} k^2 - \\delta M_\\mathswitchr {G_0}^2,\\\\\n\\hat{\\Sigma}^{\\mathswitchr {G_0}\\mathswitchr {A_0}}(k^2)&= \\Sigma^{\\mathswitchr {G_0}\\mathswitchr {A_0}}(k^2) + \\frac{1}{2}\\delta Z_{\\mathswitchr {A_0}\\mathswitchr {G_0}}(k^2-\\mathswitch {M_\\PAO}^2) +\\frac{1}{2}\\delta Z_{\\mathswitchr {G_0}\\mathswitchr {A_0}} \\;k^2 - \\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2.\n\\end{align}\n\\end{subequations}\nThe charged sector involves the following self-energies,\n\\begin{subequations}\n\\label{eq:chargedSE}\n\\begin{align}\n\\hat{\\Sigma}^{\\mathrm{H}^+ \\mathrm{H}^-}(k^2)&= \\Sigma^{\\mathrm{H}^+ \\mathrm{H}^-}(k^2) + \\delta Z_{\\mathswitchr {H^+}}(k^2-\\mathswitch {M_\\PHP}^2) - \\delta \\mathswitch {M_\\PHP}^2,\\\\\n\\hat{\\Sigma}^{\\mathrm{G}^+\\mathrm{G}^-}(k^2)&= \\Sigma^{\\mathrm{G}^+\\mathrm{G}^-}(k^2) + \\delta Z_{\\mathrm{G}^+}\\; k^2 - \\delta M_\\mathswitchr {G^+}^2,\\\\\n\\hat{\\Sigma}^{\\mathrm{G}^\\pm\\mathrm{H}^\\mp}(k^2)&= \\Sigma^{\\mathrm{G}^\\pm\\mathrm{H}^\\mp}(k^2) + \\frac{1}{2}\\delta Z_{\\mathswitchr H\\mathswitchr {G^+}}(k^2-\\mathswitch {M_\\PHP}^2) +\\frac{1}{2}\\delta Z_{\\mathswitchr {G}\\mathswitchr {H^+}}\\; k^2 - \\delta M_{\\mathswitchr {G}\\mathswitchr {H^+}}^2.\n\\end{align}\n\\end{subequations}\nThe mass mixing constants are given in Eqs.~\\eqref{eq:MHhchoice2} and \\eqref{eq:Mabchoice2}.\nOn these two-point functions we now impose our renormalization conditions. First, we fix the renormalized mass parameters to the on-shell\nvalues, so that the zeros of the real parts of the one-particle-irreducible two-point functions are located at the squares of the physical masses:\n\\begin{align}\n\\mathrm{Re}\\,\\hat{\\Sigma}^{\\mathswitchr H\\PH}(\\mathswitch {M_\\PH}^2)&=0,& \\mathrm{Re}\\,\\hat{\\Sigma}^{\\mathswitchr h\\Ph}(\\mathswitch {M_\\Ph}^2)&=0,\\nonumber\\\\\n\\mathrm{Re}\\,\\hat{\\Sigma}^{\\mathswitchr {A_0}\\PAO}(\\mathswitch {M_\\PAO}^2)&=0,&\\mathrm{Re}\\,\\hat{\\Sigma}^{\\mathrm{H}^+\\mathrm{H}^-}(\\mathswitch {M_\\PHP}^2)&=0.\n\\label{eq:massrenconddef}\n\\end{align}\nUsing Eqs.~\\eqref{eq:neutralSE}, \\eqref{eq:CPoddSE}, and \\eqref{eq:chargedSE}, fixes the mass renormalization constants to\n\\begin{align}\n\\delta \\mathswitch {M_\\PH}^2&= \\mathrm{Re}\\, \\Sigma^{\\mathswitchr H\\PH}(\\mathswitch {M_\\PH}^2),&\\delta \\mathswitch {M_\\Ph}^2&= \\mathrm{Re}\\, \\Sigma^{\\mathswitchr h\\Ph}(\\mathswitch {M_\\Ph}^2),\\nonumber\\\\\n\\delta \\mathswitch {M_\\PAO}^2&= \\mathrm{Re}\\, \\Sigma^{\\mathswitchr {A_0}\\PAO}(\\mathswitch {M_\\PAO}^2),&\\delta M_{\\mathswitchr {H^+}}^2&= \\mathrm{Re}\\, \\Sigma^{\\mathrm{H}^+\\mathrm{H}^-}(M_{\\mathswitchr H^+}^2).\n\\end{align}\nFor the propagators of the fields, we demand that the residues of the particle poles are not changed by higher-order corrections.\nThis determines the \ndiagonal field renormalization constants by conditions on the one-particle-irreducible two-point functions,\n\\begin{align}\n\\lim_{k^2\\to \\mathswitch {M_\\PH}^2}\\mathrm{Re}\\,\\frac{\\mathrm{i}\\,\\hat{\\Gamma}^{\\mathswitchr H\\PH}(k^2)}{k^2-\\mathswitch {M_\\PH}^2}&=-1,& \n\\lim_{k^2\\to \\mathswitch {M_\\Ph}^2}\\mathrm{Re}\\,\\frac{\\mathrm{i}\\,\\hat{\\Gamma}^{\\mathswitchr h\\Ph}(k^2)}{k^2-\\mathswitch {M_\\Ph}^2}&=-1,\\nonumber\\\\\n\\lim_{k^2\\to \\mathswitch {M_\\PAO}^2}\\mathrm{Re}\\,\\frac{\\mathrm{i}\\,\\hat{\\Gamma}^{\\mathswitchr {A_0}\\PAO}(k^2)}{k^2-\\mathswitch {M_\\PAO}^2}&=-1,&\n\\lim_{k^2\\to \\mathswitch {M_\\PHP}^2}\\mathrm{Re}\\,\\frac{\\mathrm{i}\\,\\hat{\\Gamma}^{\\mathrm{H}^+\\mathrm{H}^-}(k^2)}{k^2-\\mathswitch {M_\\PHP}^2}&=-1,\n\\end{align}\nwhich implies\n\\begin{align}\n\\delta Z_\\mathswitchr H&= -\\mathrm{Re}\\, \\Sigma'^{\\mathswitchr H\\PH}(\\mathswitch {M_\\PH}^2),\n&\\delta Z_\\mathswitchr h&= -\\mathrm{Re}\\, \\Sigma'^{\\mathswitchr h\\Ph}(\\mathswitch {M_\\Ph}^2), \\label{eq:dZHdef}\\nonumber\\\\\n\\delta Z_{\\mathswitchr {A_0}}&= -\\mathrm{Re}\\, \\Sigma'^{\\mathswitchr {A_0}\\PAO}(\\mathswitch {M_\\PAO}^2),\n&\\delta Z_{\\mathrm{H}^+}&= -\\mathrm{Re}\\, \\Sigma'^{\\mathrm{H}^+\\mathrm{H}^-}(\\mathswitch {M_\\PHP}^2),\n\\end{align}\nwhere we introduced \n${\\Sigma}'(k^2)$ as the derivative w.r.t.\\ the argument $k^2$.\nTo fix the mixing renormalization constants, we enforce the condition that on-mass-shell fields do not mix, i.e.\n\\begin{align}\n\\mathrm{Re}\\,\\hat{\\Sigma}^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\PH}^2)&=0,& \\mathrm{Re}\\,\\hat{\\Sigma}^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\Ph}^2)&=0,\n\\nonumber\\\\\n\\mathrm{Re}\\,\\hat{\\Sigma}^{\\mathswitchr {G_0}\\mathswitchr {A_0}}(\\mathswitch {M_\\PAO}^2)&=0,&\\mathrm{Re}\\,\\hat{\\Sigma}^{\\mathswitchr {G}^+\\mathswitchr H^-}(\\mathswitch {M_\\PHP}^2)&=0.\n\\label{eq:mixrencond}\n\\end{align}\nAfter inserting the renormalized self-energies we obtain\n\\begin{align}\n\\delta Z_{\\mathswitchr H\\mathswitchr h}&= 2 \\, \\frac{\\delta M_{\\mathswitchr H\\mathswitchr h}^2-\\mathrm{Re}\\,\\Sigma^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\Ph}^2)}{\\mathswitch {M_\\Ph}^2-\\mathswitch {M_\\PH}^2},\n&\\delta Z_{\\mathswitchr h\\mathswitchr H}&= 2 \\, \\frac{\\delta M_{\\mathswitchr H\\mathswitchr h}^2-\\mathrm{Re}\\,\\Sigma^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\PH}^2)}{\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2},\n\\nonumber\\\\\n\\delta Z_{\\mathswitchr {G_0}\\mathswitchr {A_0}}&= 2 \\, \\frac{\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2-\\mathrm{Re}\\,\\Sigma^{\\mathswitchr {G_0}\\mathswitchr {A_0}}(\\mathswitch {M_\\PAO}^2)}{\\mathswitch {M_\\PAO}^2},\n&\\delta Z_{\\mathswitchr {G}\\mathrm{H}^+}&= 2 \\, \\frac{\\delta M_{\\mathswitchr {G}\\mathswitchr {H^+}}^2-\\mathrm{Re}\\,\\Sigma^{\\mathrm{H}^+\\mathrm{G}^-}(\\mathswitch {M_\\PHP}^2)}{\\mathswitch {M_\\PHP}^2}.\\label{eq:mixfieldrendef}\n\\end{align}\nSince Goldstone-boson fields \ndo not correspond to physical states, we do not render Green functions with external \nGoldstone bosons finite, so that we need not fix the constants \n$\\delta Z_{\\mathswitchr {G_0}}$, $\\delta Z_{\\mathswitchr {A_0}\\mathswitchr {G_0}}$, $\\delta Z_{\\mathswitchr {G^+}}$, $\\delta Z_{\\mathswitchr H\\mathswitchr {G^+}}$;\nwe could even set them to zero consistently.\nThe possible \n$\\mathswitchr Z\\mathswitchr {A_0}$ and $\\mathrm{W}^\\pm \\mathswitchr H^\\mp$ mixings\nvanish for physical on-shell gauge bosons due to the Lorentz structure of the two-point function \nand the fact that polarization vectors $\\varepsilon^\\mu$ are orthogonal to the corresponding momentum.\nUsing the convention\n\\begin{align}\n\\hat\\Gamma^{\\mathswitchr Z\\mathswitchr {A_0}}_\\mu(k) &= k_\\mu \\, \\hat\\Sigma^{\\mathswitchr Z\\mathswitchr {A_0}}(k^2) = \nk_\\mu \\left[ \\Sigma^{\\mathswitchr Z\\mathswitchr {A_0}}(k^2)-\\textstyle\\frac{1}{2}\\mathswitch {M_\\PZ}\\delta Z_{\\mathswitchr {G_0}\\mathswitchr {A_0}}\\right],\n\\nonumber\\\\\n\\hat\\Gamma^{\\mathswitchr W^\\pm\\mathswitchr H^\\mp}_\\mu(k) &= k_\\mu \\, \\hat\\Sigma^{\\mathswitchr W^\\pm\\mathswitchr H^\\mp}(k^2) = \nk_\\mu \\left[ \\Sigma^{\\mathswitchr W^\\pm\\mathswitchr H^\\mp}(k^2)\\pm\\textstyle\\frac{\\mathrm{i}}{2}\\mathswitch {M_\\PW}\\delta Z_{\\mathswitchr {G}\\mathswitchr {H^+}}\\right],\n\\end{align}\nwhere all fields are incoming and $k$ is the incoming momentum of the gauge-boson fields,\nthe vector--scalar mixing self-energies obey\n\\begin{align}\n\\left.\\varepsilon^\\mu_\\mathswitchr Z k_\\mu \\,\\mathop{\\mathrm{Re}}\\nolimits\\,\\hat{\\Sigma}^{\\mathswitchr Z\\mathswitchr {A_0}}(k^2)\\right|_{k^2= \\mathswitch {M_\\PZ}^2}=0,\n\\qquad &\n\\left.\\varepsilon^\\mu_\\mathswitchr W k_\\mu \\,\\mathop{\\mathrm{Re}}\\nolimits\\,\\hat{\\Sigma}^{\\mathswitchr W^\\pm\\mathswitchr H^\\mp}(k^2)\\right|_{k^2= \\mathswitch {M_\\PW}^2}=0.\n\\end{align}\nThe mixing self-energies on the other on-shell points \n$k^2=\\mathswitch {M_\\PAO}^2$ and $k^2=\\mathswitch {M_\\PHP}^2$, respectively, \nare connected to the mixing of \n$\\mathswitchr {A_0}$ or $\\mathswitchr H^\\pm$ with the Goldstone-boson fields of the\n$\\mathswitchr Z$ or the $\\mathswitchr W$ boson and can be calculated from a BRST symmetry \\cite{Becchi:1975nq}. The BRST variation of the Green functions of one anti-ghost and a Higgs field\n\\begin{align}\n \\delta_\\mathrm{BRST} \\langle 0 | T \\bar{u}^\\mathswitchr Z(x) A_0(y)|0\\rangle=0, &&\n \\delta_\\mathrm{BRST} \\langle 0 | T \\bar{u}^\\pm(x) H^\\pm(y)|0\\rangle=0,\n\\end{align}\nimplies Slavnov--Taylor identities. While the variation of the anti-ghost fields yields the gauge-fixing term, the variation of the Higgs fields introduces ghost contributions which \nvanish for on-shell momentum resulting in%\n\\footnote{A particularly simple, alternative \nway to derive these identities is to exploit the gauge invariance\nof the effective action in the background-field gauge, as \ndone in \\citere{Denner:1994xt} for the SM. The respective Ward identities for the background\nfields differ from the Slavnov--Taylor identities only by off-shell terms, which vanish\non the particle poles. Generalizing the derivation of \\citere{Denner:1994xt} to the THDM\nand adapting the results to our conventions for self-energies,\nthe desired Ward identities for the unrenormalized background fields read\n\\begin{align}\n0 &= k^2 {\\Sigma}^{\\hat\\mathswitchr Z\\hat{\\mathrm{PA}}_0}(k^2)+ \\mathswitch {M_\\PZ} {\\Sigma}^{\\hat{\\mathrm{PA}}_0\\hat\\mathswitchr {G}_0}(k^2)\n+\\frac{e}{2\\mathswitch {c_{\\scrs\\PW}}\\mathswitch {s_{\\scrs\\PW}}}\\left(T^{\\hat\\mathswitchr H}s_{\\beta-\\alpha}-T^{\\hat\\mathswitchr h}c_{\\beta-\\alpha}\\right),\n\\\\\n0 &= \nk^2 {\\Sigma}^{\\hat\\mathswitchr W^\\pm\\hat\\mathswitchr H^\\mp}(k^2)\\mp \\mathrm{i} \\mathswitch {M_\\PW} {\\Sigma}^{\\hat\\mathswitchr {G}^\\pm\\hat\\mathswitchr H^\\mp}(k^2)\n\\mp\\frac{\\mathrm{i} e}{2\\mathswitch {s_{\\scrs\\PW}}}\\left(T^{\\hat\\mathswitchr H}s_{\\beta-\\alpha}-T^{\\hat\\mathswitchr h}c_{\\beta-\\alpha}\\right),\n\\end{align}\nwhere the carets on the fields indicate background fields.\nSetting $k^2$ to $\\mathswitch {M_\\PAO}^2$ or $\\mathswitch {M_\\PHP}^2$, respectively, and adding the \nrelevant renormalization constants, directly leads to the identities \\refeq{eq:ZA-WI}\nand \\refeq{eq:WH-WI}.\n}\n\\begin{align}\n\\label{eq:ZA-WI}\n \\left[k^2 \\hat{\\Sigma}^{\\mathswitchr Z\\mathswitchr {A_0}}(k^2)+ \\mathswitch {M_\\PZ} \\hat{\\Sigma}^{\\mathswitchr {G_0}\\mathswitchr {A_0}}(k^2)\\right]_{k^2= \\mathswitch {M_\\PAO}^2}=0,\\\\\n \\left[k^2 \\hat{\\Sigma}^{\\mathrm{W}^\\pm\\mathswitchr H^\\mp}(k^2)\\mp \\mathrm{i} \\mathswitch {M_\\PW} \\hat{\\Sigma}^{\\mathrm{G}^\\pm\\mathrm{H}^\\mp}(k^2)\\right]_{k^2= \\mathswitch {M_\\PHP}^2}=0.\n\\label{eq:WH-WI}\n\\end{align}\nWe have verified these identities analytically and numerically.\nTogether with the renormalization condition of \nEq.~\\eqref{eq:mixrencond} we conclude that \n\\begin{align}\n\\hat{\\Sigma}^{\\mathswitchr Z\\mathswitchr {A_0}}(\\mathswitch {M_\\PAO}^2)=0, \\quad \\hat{\\Sigma}^{\\mathrm{W}^\\pm\\mathswitchr H^\\mp}(\\mathswitch {M_\\PHP}^2)=0.\n\\end{align}\nThis set of renormalization conditions ensures that no on-shell two-point vertex function obtains any one-loop corrections,\nand the corresponding external\nself-energy diagrams do not have to be taken into account in any calculation. \n\n\\subsubsection{Electroweak sector}\n\nThe fixing of the renormalization constants of the electroweak sector is identical to the SM case. The mass renormalization constants are fixed in such a way that the squares of the masses correspond to the \n(real parts of the) locations of the poles of the gauge-boson propagators. The field renormalization constants are fixed by the conditions that residues of on-shell gauge-boson propagators do not obtain higher-order corrections, and that on-shell gauge bosons do not mix. For a better bookkeeping we also keep the dependent renormalization constants $\\delta \\mathswitch {c_{\\scrs\\PW}}$ and $\\delta \\mathswitch {s_{\\scrs\\PW}}$ in our calculation. This results in~\\cite{Denner:1991kt}\n\\begin{align}\n\\delta \\mathswitch {M_\\PW}^2 &=\\mathrm{Re}\\, \\Sigma^W_\\mathrm{T}(\\mathswitch {M_\\PW}^2),&\\delta Z_{\\mathswitchr W}&=-\\mathrm{Re}\\, \\Sigma'^{\\mathswitchr W}_\\mathrm{T}(\\mathswitch {M_\\PW}^2),\\nonumber\\\\\n\\delta \\mathswitch {M_\\PZ}^2&=\\mathrm{Re}\\, \\Sigma^{\\mathswitchr Z\\PZ}_\\mathrm{T}(\\mathswitch {M_\\PZ}^2),\\nonumber\\\\\n\\delta Z_{\\mathswitchr Z\\PZ}&=-\\mathrm{Re}\\, \\Sigma'^{\\mathswitchr Z\\PZ}_\\mathrm{T}(\\mathswitch {M_\\PZ}^2),&\\delta Z_{{\\mathrm{PA}}\\PA}&=-\\mathrm{Re}\\, \\Sigma'^{{\\mathrm{PA}}\\PA}_\\mathrm{T}(0),\\nonumber\\\\\n\\delta Z_{{\\mathrm{PA}}\\mathswitchr Z}&=-2\\mathrm{Re}\\, \\frac{\\Sigma^{{\\mathrm{PA}}\\mathswitchr Z}_\\mathrm{T}(\\mathswitch {M_\\PZ}^2)}{\\mathswitch {M_\\PZ}^2},&\n\\delta Z_{\\mathswitchr Z{\\mathrm{PA}}}&=2\\mathrm{Re}\\, \\frac{\\Sigma^{{\\mathrm{PA}}\\mathswitchr Z}_\\mathrm{T}(0)}{\\mathswitch {M_\\PZ}^2},\\nonumber\\\\\n\\delta \\mathswitch {c_{\\scrs\\PW}}&=\\frac{\\mathswitch {c_{\\scrs\\PW}}}{2}\\left(\\frac{\\delta \\mathswitch {M_\\PW}^2}{\\mathswitch {M_\\PW}^2}-\\frac{\\delta \\mathswitch {M_\\PZ}^2}{\\mathswitch {M_\\PZ}^2}\\right), &\\delta \\mathswitch {s_{\\scrs\\PW}}&= -\\frac{\\mathswitch {c_{\\scrs\\PW}}}{\\mathswitch {s_{\\scrs\\PW}}}\\delta \\mathswitch {c_{\\scrs\\PW}}.\n\\end{align}\nThe electric charge $e$ is defined via the $\\mathrm{ee} \\gamma $ coupling in the Thomson limit of on-shell external electrons and zero momentum transfer to the photon,\nwhich yields in BHS convention \\cite{Denner:1991kt}\n\\begin{align}\n\\delta Z_e = -\\frac{1}{2} \\left(\\delta Z_{{\\mathrm{PA}}\\PA}+ \\frac{\\mathswitch {s_{\\scrs\\PW}}}{\\mathswitch {c_{\\scrs\\PW}}} \\delta Z_{\\mathswitchr Z{\\mathrm{PA}}}\\right).\n\\label{eq:dZe}\n\\end{align}\n\n\\subsubsection{Fermions}\n\nThe renormalization conditions for the fermions are identical to the ones in the SM, described in detail in Ref.~\\cite{Denner:1991kt}. We demand that the (real parts of the locations of the)\npoles of the fermion propagators correspond to the squared fermion masses, \n and that on-shell fermion propagators do not obtain loop corrections. \nAssuming the CKM matrix \nequal to the unit matrix, the results for the renormalization constants simplify to\n\\begin{align}\n\\delta m_{f,i}&=\\frac{m_{f,i}}{2}\\mathrm{Re}\\,\\big[\\Sigma^{f,{\\mathrm{L}}}_i(m^2_{f,i})+\\Sigma^{f,{\\mathrm{R}}}_i(m^2_{f,i})+2 \\Sigma^{f,{\\mathrm{S}}}_i(m^2_{f,i})\\big],\n\\nonumber\\\\\n\\delta Z^{f,\\sigma}_i&=-\\mathrm{Re}\\, \\Sigma^{f,\\sigma}_i(m^2_{f,i})-m_{f,i}^2\\frac{\\partial}{\\partial k^2} \\mathrm{Re} \\big[\\Sigma_i^{f,{\\mathrm{L}}}(k^2)+\\Sigma^{f,{\\mathrm{R}}}_i(k^2)+2\\Sigma^{f,{\\mathrm{S}}}_i(k^2)\\big]\\Big|_{k^2=m_{f,i}^2}, \\quad \\sigma={\\mathrm{L}},{\\mathrm{R}},\n\\end{align}\nwhere we have \nused the usual \ndecomposition of the fermion self-energies into a left-handed, a right-handed, and a scalar part, $\\Sigma^{f,{\\mathrm{L}}}_i$, $\\Sigma^{f,{\\mathrm{R}}}_i$, and $ \\Sigma^{f,{\\mathrm{S}}}_i$, respectively.\nThe expressions for a non-trivial CKM matrix can be found in \\citere{Denner:1991kt}.\n\n\\subsection{\\boldmath{$\\overline{\\mathrm{MS}}$ renormalization conditions}}\n\\label{sec:diffrenschemes}\n\nIn the four renormalization schemes we \nare going to present, the imposed on-shell conditions are identical, and the differences only occur in the choice of different $\\overline{\\mathrm{MS}}$ conditions.\nThe parameters $\\alpha$ or $\\lambda_3$ governing the mixing of the CP-even Higgs bosons, and the parameters $\\beta$ and $\\lambda_5$ need to be fixed.\nA formulation of an on-shell condition for these parameters is not obvious. \nOne could relate the parameters to \nsome physical processes, such as Higgs-boson decays, and demand that these processes do not \nreceive higher-order corrections. However, so far, no sign of further Higgs bosons has been observed, hence, there is no distinguished process, and such a prescription does not only require more calculational effort, but could introduce artificially large corrections to the corresponding parameters, which would spread to many other observables, as discussed in Refs.~\\cite{Freitas:2002um,Krause:2016oke}. Therefore, we choose to renormalize these parameters within the $\\overline{\\mathrm{MS}}$ scheme,\nthough different variables (such as $\\alpha$ or $\\lambda_3$) can be chosen to parameterize the model. Imposing an $\\overline{\\mathrm{MS}}$ condition on either of the parameters leads to differences in the calculation of observables. In addition, gauge-dependent definitions of $\\overline{\\mathrm{MS}}$-renormalized parameters spoil the gauge independence of the relations between input parameters and observables. However, gauge dependences might be even\nacceptable if the renormalization scheme yields stable results and a good convergence of the perturbation series. The price to pay is that subsequent calculations should be done in the same gauge\nor properly translated into another gauge. We will discuss different renormalization schemes based on different treatments of\n$\\alpha$ or $\\lambda_3$ parameterizing the CP-even Higgs-boson mixing, of the parameter $\\beta$, and of the Higgs coupling constant\n $\\lambda_5$ \nin the following. We begin with the so-called $\\overline{\\mathrm{MS}}(\\alpha)$ scheme.\n\n\\subsubsection{\\boldmath{$\\overline{\\mathrm{MS}}(\\alpha)$ scheme}}\n\\label{sec:alphaMSscheme}\nIn this scheme the independent parameter set is $\\{p_{\\mathrm{mass}}\\}$ of Eq.~\\eqref{eq:physparaNLO}, so that the parameters $\\beta$, $\\alpha$, and $\\lambda_5$ are renormalized in $\\overline{\\mathrm{MS}}$. The corresponding counterterm Lagrangian was derived in Sect.~\\ref{sec:renormassb}.\n\n\\paragraph{The renormalization constant \\boldmath{$\\delta \\beta$}:}\n\nThe renormalization constant \n$\\delta \\tan{\\beta}=\\delta \\beta\/{c_\\beta^2}$\nof the mixing angle $\\beta$ is\nrelated to the renormalization constants of the vevs\nby demanding the defining relation $\\tan\\beta=v_2\/v_1$ for bare and renormalized \nquantities.\nIn {$\\overline{\\mathrm{MS}}$}{}, $\\delta \\beta$ can be most easily calculated using the minimal field\nrenormalization \\refeq{eq:inputfieldren} with the following renormalization transformation\nof the vevs,\n\\begin{align}\nv_{1,0} &= Z_{\\mathswitchr H_1}^{1\/2}(v_1+\\delta \\overline{v}_1),&v_{2,0}&=Z_{\\mathswitchr H_2}^{1\/2}(v_2+\\delta \\overline{v}_2),\n \\label{eq:vevren}\n\\end{align}\nUsing the well-known relation~\\cite{Sperling:2013eva}\n\\begin{align}\n\\delta \\overline{v}_1\/v_1-\\delta \\overline{v}_2\/v_2=\\mathrm{finite},\\label{eq:vacrencond}\n\\end{align}\nthe general form of $\\delta\\beta$ in the {$\\overline{\\mathrm{MS}}$}{} scheme\n\\begin{align}\n\\delta \\beta\n&= \\frac{s_{2\\beta}}{2}\\left(\n\\frac{\\delta \\overline{v}_2}{v_2} -\\frac{\\delta \\overline{v}_1}{v_1}- \\frac{1}{2}\\delta Z_{\\mathswitchr H_1} + \\frac{1}{2} \\delta Z_{\\mathswitchr H_2} \n\\right)\\bigg|_\\mathrm{UV}\n\\end{align}\nsimplifies to \n\\begin{align}\n\\delta \\beta\n&=\\frac{s_{2\\beta}}{4}(-\\delta Z_{\\mathswitchr H_1}+\\delta Z_{\\mathswitchr H_2})\\big|_\\mathrm{UV}\n=\\frac{s_{2\\beta}}{4c_{2\\alpha}} \\left(\\delta Z_\\mathswitchr h -\\delta Z_\\mathswitchr H\\right)\\big|_\\mathrm{UV}\n=\\frac{s_{2\\beta}}{4s_{2\\alpha}} \\left(\\delta Z_{\\mathswitchr h\\mathswitchr H} +\\delta Z_{\\mathswitchr H\\mathswitchr h}\\right)\\big|_\\mathrm{UV},\n\\label{eq:tanbetarencond}\n\\end{align}\nwhere $\\big|_\\mathrm{UV}$ indicates that we take only the UV-divergent parts, which are proportional to \nthe standard divergence $\\Delta_\\mathrm{UV}$.\nThe explicit calculation of the UV-divergent terms of $\\delta Z_\\mathswitchr h,$ $\\delta Z_\\mathswitchr H$ according to Eqs.~\\eqref{eq:dZHdef} in \n't~Hooft-Feynman gauge reveals that only diagrams with closed fermion loops contribute to the counterterm, \n\\begin{align}\\label{eq:deltabeta}\n\\delta \\beta = -\\Delta_\\mathrm{UV}\\, \\frac{e^2}{64\\pi^2\\mathswitch {M_\\PW}^2\\mathswitch {s_{\\scrs\\PW}}^2} \\sum_f c_f \\xi^f_\\mathswitchr {A_0} m_f^2,\n\\end{align}\nwith the colour factors $c_\\mathrm{quark}=3$, $c_\\mathrm{lepton}=1$ and the coupling coefficients $\\xi^f_{\\mathswitchr {A_0}}$ \nas defined in Tab.~\\ref{tab:yukint}. In the class of $R_\\xi$ gauges this result is gauge independent at one-loop order \\cite{Krause:2016oke,Denner:2016etu}.\n\n\\paragraph{Neutral Higgs mixing:}\nIn the neutral Higgs sector, relations between field renormalization constants can also be used to determine another parameter in $\\overline{\\mathrm{MS}}$.\nThe first four equations of Eqs.~\\eqref{eq:fieldrenconstrel} can be solved for $\\delta \\alpha$ in various\nways, e.g., yielding\n\\begin{align}\n\\label{eq:deltaalphaderiv}\n\\delta \\alpha\\big|_\\mathrm{UV}&= \\frac{1}{4} \\left(\\delta Z_{\\mathswitchr H\\mathswitchr h}-\\delta Z_{\\mathswitchr h\\mathswitchr H}\\right) \\big|_\\mathrm{UV}.\n\\end{align}\nThe field renormalization constants can be inserted according to Eqs.~\\eqref{eq:dZHdef}, \\eqref{eq:mixfieldrendef} using \n$\\delta M_{\\mathswitchr H\\mathswitchr h}^2 = 0$, thus \n\\begin{align}\n\\label{eq:deltaalphaderiv2}\n\\delta \\alpha&= \\left. \\mathrm{Re}\\frac{\\Sigma^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\PH}^2)+\\Sigma^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\Ph}^2)}{2(\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2)}\\right|_\\mathrm{UV}.\n\\end{align}\nAn explicit calculation of the counterterm \nyields for the fermionic contribution, \n\\begin{align}\\label{eq:delalphaferm}\n\\delta \\alpha\\big|_\\mathrm{ferm} = \\Delta_\\mathrm{UV}\\, \\frac{e^2 s_{2\\alpha}}{64\\pi^2\\mathswitch {M_\\PW}^2\\mathswitch {s_{\\scrs\\PW}}^2 s_{2\\beta}(\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2)} \n\\sum_f c_f \\xi^f_\\mathswitchr {A_0} m_f^2 (\\mathswitch {M_\\PH}^2+\\mathswitch {M_\\Ph}^2-12m_f^2),\n\\end{align}\nand for the bosonic contribution\n\\begin{align}\\label{eq:delalphabos}\n\\delta \\alpha\\big|_\\mathrm{bos} = &\n-\\Delta_\\mathrm{UV} \\frac{\\lambda_5^2 \\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2 }{8\\pi^2 e^2 (\\mathswitch {M_\\Ph}^2 - \\mathswitch {M_\\PH}^2) s_{2\\beta}^2}\n \\Big[s_{2 (\\alpha - 3 \\beta)} + 10 s_{2 (\\alpha - \\beta)} + 13 s_{2 (\\alpha + \\beta)}\\Big]\n\\nonumber\\\\\n& + \\Delta_\\mathrm{UV} \\frac{\\lambda_5 }{ 128\\pi^2 (\\mathswitch {M_\\Ph}^2 - \\mathswitch {M_\\PH}^2) s_{2\\beta}^2} \\Big[\n -4 \\mathswitch {M_\\PH}^2 (13 c_{2 \\alpha} + 2 c_{2 (\\alpha - 2 \\beta)} - 27 c_{2 \\beta}) s_{2 \\alpha} \n\\nonumber\\\\\n& + 4 \\mathswitch {M_\\Ph}^2 (13 c_{2 \\alpha} + 2 c_{2 (\\alpha - 2 \\beta)} + 27 c_{2 \\beta}) s_{2 \\alpha} \n + 2 \\mathswitch {M_\\PHP}^2 (s_{2 (\\alpha - 3 \\beta)} - 6 s_{2 (\\alpha - \\beta)} + 13 s_{2 (\\alpha + \\beta)}) \n\\nonumber\\\\\n& - \\mathswitch {M_\\PAO}^2 (7 s_{2 (\\alpha - 3 \\beta)} + 86 s_{2 (\\alpha - \\beta)} + 91 s_{2 (\\alpha + \\beta)})\n + 4 (2 \\mathswitch {M_\\PW}^2 + \\mathswitch {M_\\PZ}^2) s_{2 (\\alpha - \\beta)} s_{2 \\beta}^2 \n\\Big]\n\\nonumber\\\\\n& + \\Delta_\\mathrm{UV} \\frac{e^2 }{ 1024\\pi^2 (\\mathswitch {M_\\Ph}^2 - \\mathswitch {M_\\PH}^2) \\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2 s_{2\\beta}^2} \\Big[\n- 2 \\mathswitch {M_\\PH}^4 (-36 c_{2 \\alpha} + 5 c_{4 \\alpha - 2 \\beta} + 31 c_{2 \\beta}) s_{2 \\alpha} \n\\nonumber\\\\\n& + 4 \\mathswitch {M_\\Ph}^2 \\mathswitch {M_\\PH}^2 (5 c_{4 \\alpha - 2 \\beta} - 29 c_{2 \\beta}) s_{2 \\alpha} \n- 2 \\mathswitch {M_\\Ph}^4 (36 c_{2 \\alpha} + 5 c_{4 \\alpha - 2 \\beta} + 31 c_{2 \\beta}) s_{2 \\alpha} \n\\nonumber\\\\\n& + 32 \\mathswitch {M_\\PHP}^4 s_{2 (\\alpha - \\beta)} s_{2 \\beta}^2 \n+ 2 \\mathswitch {M_\\PH}^2 \\mathswitch {M_\\PHP}^2 (3 s_{4 \\alpha} + 4 s_{2 (\\alpha - \\beta)} + s_{4 (\\alpha - \\beta)} \n\t+ 9 s_{4 \\beta} - 12 s_{2 (\\alpha + \\beta)}) \n\\nonumber\\\\\n& - 2 \\mathswitch {M_\\Ph}^2 \\mathswitch {M_\\PHP}^2 (3 s_{4 \\alpha} -4 s_{2 (\\alpha - \\beta)} + s_{4 (\\alpha - \\beta)} \n\t+ 9 s_{4 \\beta} + 12 s_{2 (\\alpha + \\beta)}) \n\\nonumber\\\\\n& - 2 \\mathswitch {M_\\PAO}^4 (5 s_{2 (\\alpha - 3 \\beta)} + 42 s_{2 (\\alpha - \\beta)} + 41 s_{2 (\\alpha + \\beta)}) \n\\nonumber\\\\\n& + \\mathswitch {M_\\PAO}^2 \\mathswitch {M_\\PH}^2 (-49 s_{4 \\alpha} + 112 s_{2 (\\alpha - \\beta)} - 7 s_{4 (\\alpha - \\beta)} \n\t+ s_{4 \\beta} + 96 s_{2 (\\alpha + \\beta)})\n\\nonumber\\\\\n& + \\mathswitch {M_\\PAO}^2 \\mathswitch {M_\\Ph}^2 (49 s_{4 \\alpha} + 112 s_{2 (\\alpha - \\beta)} + 7 s_{4 (\\alpha - \\beta)} \n\t- s_{4 \\beta} + 96 s_{2 (\\alpha + \\beta)}) \n\\nonumber\\\\\n& + 4 \\mathswitch {M_\\PAO}^2 \\mathswitch {M_\\PHP}^2 (s_{2 (\\alpha - 3 \\beta)} - 6 s_{2 (\\alpha - \\beta)} + 13 s_{2 (\\alpha + \\beta)}) \n\\nonumber\\\\\n& + 4 (2 \\mathswitch {M_\\PW}^2 + \\mathswitch {M_\\PZ}^2) s_{2 (\\alpha - \\beta)} s_{2 \\beta} \n\t((\\mathswitch {M_\\Ph}^2 - \\mathswitch {M_\\PH}^2) s_{2 \\alpha} + 2 \\mathswitch {M_\\PAO}^2 s_{2 \\beta}) \n\\nonumber\\\\\n& + 48 (2 \\mathswitch {M_\\PW}^4 + \\mathswitch {M_\\PZ}^4) s_{2 (\\alpha - \\beta)} s_{2 \\beta}^2 \n\\Big].\n\\end{align}\nThis result, which is derived in 't~Hooft Feynman gauge, is gauge dependent \\cite{Krause:2016oke,Denner:2016etu}.\n\n\\paragraph{Higgs self-coupling:}\nThe Higgs self-coupling counterterm $\\delta \\lambda_5$ has to be fixed via a vertex correction. We define this renormalization constant in $\\overline{\\mathrm{MS}}$ as well, as there is no distinguished process to fix it on-shell. Any 3- or 4-point vertex function with external Higgs bosons is suited to calculate the divergent terms.\nSince the $\\mathswitchr H\\mathswitchr {A_0}\\PAO$ vertex correction involves fewest diagrams, it is our preferred choice. The condition is\n\\begin{align}\n\\hat{\\Gamma}^{\\mathswitchr H\\mathswitchr {A_0}\\PAO}\\big|_\\mathrm{UV}=\n\\parbox{2.4cm}{\\includegraphics[scale=1]{.\/diagrams\/L5ren.eps}}\\left.\\rule{0cm}{1.1cm}\\right|_\\mathrm{UV}=0.\n\\end{align}\nSolving this equation for $\\delta \\lambda_5$ fixes this renormalization constant. The generic one-loop diagrams appearing in this vertex correction are shown in Fig.~\\ref{fig:dL5diags}, the contribution of the diagrams involving closed fermion loops is \n\\begin{align}\n\\label{eq:L5ferm}\n\\delta \\lambda_{5,\\mathrm{ferm}}=\\Delta_\\mathrm{UV} \n\\frac{e^2 \\lambda_5}{16\\pi^2 \\mathswitch {M_\\PW}^2\\mathswitch {s_{\\scrs\\PW}}^2}\\sum_f c_f \\left(1+\\frac{c_{2\\beta}}{s_{2\\beta}}\\xi_{\\mathswitchr {A_0}}^f\\right) m_f^2.\n\\end{align}\n\\begin{figure}\n\\vspace{-15pt}\n\\input{diagrams\/triangles}\n\\vspace{-20pt}\n\\caption{Generic diagrams contributing to the $HA_0A_0$ vertex correction used for the renormalization of $\\lambda_5$.}\n\\label{fig:dL5diags}\n\\end{figure}\nThe diagrams containing only bosons lead to \n\\begin{align}\n\\label{eq:L5bos}\n\\delta \\lambda_5\\big|_\\mathrm{bos}{}={}& {}\n\\Delta_\\mathrm{UV} \\frac{\\lambda_5}{32 \\pi^2} \n\\left(2 \\lambda_1 + 2\\lambda_2+8\\lambda_3 +12 \\lambda_4- 9 g_2^2 - 3 g_1^2\\right)\n\\nonumber\\\\\n{}={}& - \\Delta_\\mathrm{UV}\\frac{\\lambda_5^2c^2_{2\\beta}}{4 \\pi^2 s^2_{2\\beta}}\n +\\Delta_\\mathrm{UV}\\frac{\\lambda_5 e^2}{64 \\pi^2\\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2 s_{2\\beta}^2} \\Big[\n \\mathswitch {M_\\PH}^2 (2 + c_{2 (\\alpha - \\beta)} - 3 c_{2 (\\alpha + \\beta)}) \n\\nonumber\\\\\n& + \\mathswitch {M_\\Ph}^2 (2 - c_{2 (\\alpha - \\beta)} + 3 c_{2 (\\alpha + \\beta)})\n+ \\mathswitch {M_\\PAO}^2 (1 - 5 c_{4 \\beta})\n- 4 \\mathswitch {M_\\PHP}^2 s_{2 \\beta}^2\n- 6(2 \\mathswitch {M_\\PW}^2+\\mathswitch {M_\\PZ}^2) s_{2 \\beta}^2 \n\\Big].\n\\end{align}\nSince $\\lambda_5$ is a fundamental parameter of the Higgs potential, an $\\overline{\\mathrm{MS}}$ definition leads to a gauge-independent counterterm.\n\n\\subsubsection{\\boldmath{$\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme}}\n\\label{sec:lambda3MSscheme}\n\nIn this scheme, the independent parameter set is $\\{p'_{\\mathrm{mass}}\\}$ defined in Eq.~\\eqref{eq:physparaNLOprime}. The renormalization of $\\beta$ and $\\lambda_5$ is identical to the previous renormalization scheme and not stated again, but now the parameter $\\lambda_3$ (instead of $\\alpha$) is an independent parameter being renormalized in $\\overline{\\mathrm{MS}}$. This has the advantage that this parameter is gauge independent, as it is a defining parameter of the basic \nparameterization of the Higgs potential\nand thus is safe against potentially gauge-dependent contributions appearing in relations between bare parameters. \nAs stated above, the {$\\overline{\\mathrm{MS}}$}{} renormalization of the \nparameter $\\beta$ generally breaks gauge \nindependence, but in $\\mathrm{R}_\\xi$ gauges the gauge dependence cancels at one loop \\cite{Krause:2016oke,Denner:2016etu}, \nso that this scheme yields gauge-independent results at NLO. \nWe take the counterterm potential of Sect.~\\ref{sec:renormassb}, but treat $\\delta \\alpha$ as a dependent counterterm. \nAs $\\alpha$ is a pure mixing angle, we choose to apply the renormalization prescription of Sect.~\\ref{sec:renormassb}, where the mixing angle diagonalizes the potential to all orders.\nThe relation between $\\delta \\alpha$ and the independent constants is given in Eq.~\\eqref{eq:NLOmassmixHh} with $ \\delta M_{\\mathswitchr H\\mathswitchr h}^2=0$,\n\\begin{align}\n\\label{eq:alphadef32}\n \\delta \\alpha=\\frac{f_\\alpha(\\{ \\delta p'_\\mathrm{mass}\\})} {\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2},\n\\end{align}\nwhere $f_\\alpha(\\{ \\delta p'_\\mathrm{mass}\\})$ can be obtained from Eq.~\\eqref{eq:Hhmassmixing} by applying the renormalization transformation of Eq.~\\eqref{eq:physpararen1} (which is identical to the renormalization transformation of Sect.~\\ref{sec:renormassb}, but renormalizing $\\lambda_3$ instead of $\\alpha$). This yields\n\\begin{align}\nf_\\alpha(\\{\\delta p'_\\mathrm{mass}\\})={}&\n\\frac{1}{2} t_{2 \\alpha } \\left(\\delta \\mathswitch {M_\\Ph}^2-\\delta \\mathswitch {M_\\PH}^2\\right)\n+\\frac{s_{2 \\beta } \\left(\\delta \\mathswitch {M_\\PAO}^2-2 \\delta \\mathswitch {M_\\PHP}^2\\right)}{2 c_{2 \\alpha }}\n\\nonumber\\\\\n&+\\frac{\\delta \\beta c_{2 \\beta } \\left(\\mathswitch {M_\\PH}^2- \\mathswitch {M_\\Ph}^2\\right) t_{2 \\alpha }}{s_{2 \\beta }}\n+\\frac{2 \\mathswitch {M_\\PW}^2 s_{2\\beta } (\\delta \\lambda_3+\\delta \\lambda_5) \\mathswitch {s_{\\scrs\\PW}}^2}{e^2 c_{2 \\alpha }}\n\\nonumber\\\\\n&+\\frac{s_{2 \\beta } \\left(\\mathswitch {M_\\PAO}^2-2 \\mathswitch {M_\\PHP}^2\\right)+(\\mathswitch {M_\\Ph}^2-\\mathswitch {M_\\PH}^2)s_{2\\alpha}}{c_{2\\alpha}}\n\\left( \\delta Z_e -\\frac{\\delta \\mathswitch {s_{\\scrs\\PW}}}{\\mathswitch {s_{\\scrs\\PW}}} -\\frac{\\delta \\mathswitch {M_\\PW}^2}{2\\mathswitch {M_\\PW}^2} \\right)\n\\nonumber\\\\\n&\n-\\frac{ e \\left[\n\\delta t_\\mathswitchr H \\left(s_{\\alpha -3 \\beta }+3 s_{\\alpha +\\beta }\\right)\n+\\delta t_\\mathswitchr h \\left(c_{\\alpha -3 \\beta }+3 c_{\\alpha +\\beta } \\right) \n\\right]}{8 \\mathswitch {M_\\PW}\n c_{2 \\alpha } \\mathswitch {s_{\\scrs\\PW}}}.\\label{eq:MHhdependence}\n\\end{align}\nThe UV-divergent term of $\\delta \\alpha$ has been calculated in Eq.~\\eqref{eq:deltaalphaderiv2}, and by renormalizing $\\delta \\lambda_3$ in $\\overline{\\mathrm{MS}}$ scheme, it is clear that the dependent $\\delta \\alpha$ must now have a finite part in addition. \nWe choose this finite term in such a way that the finite part in $\\delta \\lambda_3$ \n(which results from $\\delta \\lambda_3$ by setting $\\Delta_{\\mathrm{UV}}$ to zero) vanishes and obtain\n\\begin{align}\n\\label{eq:alphainL3MS}\n\\delta \\alpha\\big|_{\\overline{\\mathrm{MS}}(\\lambda_3)} =\n\\left. \\mathrm{Re}\\frac{\\Sigma^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\PH}^2)+\\Sigma^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\Ph}^2)}{2(\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2)}\\right|_\\mathrm{UV}\n+\\frac{f_\\alpha(\\{ \\delta p'_\\mathrm{mass}\\})} {\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2}\\bigg|_\\mathrm{finite},\n\\end{align}\nwhere $\\delta \\lambda_3$ drops out as it has no finite part. The divergent part of $\\delta \\lambda_3$ can be calculated by solving Eq.~\\eqref{eq:alphadef32} and using the knowledge about the divergent parts of $\\delta \\alpha$ from Eqs.~\\eqref{eq:delalphaferm} and \\eqref{eq:delalphabos}. This results in \n\\begin{align}\n\\delta \\lambda_3= &\\biggl[\n\\frac{e^2 c_{2\\alpha}}{4 \\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2 s_{2\\beta}} \n\\left(\\mathrm{Re}\\,\\Sigma^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\PH}^2)+\\mathrm{Re}\\,\\Sigma^{\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\Ph}^2)\\right)\n-\\frac{\\delta \\beta e^2 s_{2 \\alpha} c_{2 \\beta} \\left(\\mathswitch {M_\\PH}^2- \\mathswitch {M_\\Ph}^2\\right) }{2 \\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2 s_{2 \\beta }^2}\n\\nonumber\\\\\n& -\\delta \\lambda_5\n-\\frac{e^2 s_{2 \\alpha }}{4 \\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2 s_{2\\beta}} \\left(\\delta \\mathswitch {M_\\Ph}^2-\\delta \\mathswitch {M_\\PH}^2\\right)\n-\\frac{e^2 \\left(\\delta \\mathswitch {M_\\PAO}^2-2 \\delta \\mathswitch {M_\\PHP}^2\\right)}{4 \\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2}\n\\nonumber\\\\\n&-\\frac{e^2\\left(s_{2\\beta}\\left(\\mathswitch {M_\\PAO}^2-2 \\mathswitch {M_\\PHP}^2\\right)+(\\mathswitch {M_\\Ph}^2-\\mathswitch {M_\\PH}^2) s_{2 \\alpha }\\right)}{2\\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2 s_{2\\beta}}\n\\left(\\delta Z_e- \\frac{\\delta \\mathswitch {s_{\\scrs\\PW}}}{\\mathswitch {s_{\\scrs\\PW}}} -\\frac{\\delta \\mathswitch {M_\\PW}^2}{2 \\mathswitch {M_\\PW}^2} \\right)\n\\nonumber\\\\\n& +\\frac{ e^3 \\left[\n\\delta t_\\mathswitchr H \\left(s_{\\alpha -3 \\beta }+3 s_{\\alpha +\\beta }\\right)\n+\\delta t_\\mathswitchr h \\left(c_{\\alpha -3 \\beta }+3 c_{\\alpha +\\beta } \\right) \n\\right]}{16 \\mathswitch {M_\\PW}^3 \\mathswitch {s_{\\scrs\\PW}}^3 s_{2\\beta}}\\biggr]_{\\mathrm{UV}}.\n\\label{eq:deltaL3}\n\\end{align}\nThe fermionic contribution to $\\delta \\lambda_3$ is given by\n\\begin{align}\n\\nonumber\n\\delta \\lambda_3\\big|_\\mathrm{UV,ferm}\n& = -\\delta \\lambda_5\\big|_\\mathrm{UV,ferm}\n-\\Delta_\\mathrm{UV} \\frac{3e^4}{32\\pi^2\\mathswitch {M_\\PW}^4\\mathswitch {s_{\\scrs\\PW}}^4} \\sum_i\n\\left(\\xi_\\mathswitchr {A_0}^u-\\xi_\\mathswitchr {A_0}^d\\right)^2\nm_{u,i}^2 m_{d,i}^2\n\\\\ \n&\\quad\n- \\Delta_\\mathrm{UV} \\frac{e^4}{64\\pi^2\\mathswitch {M_\\PW}^4\\mathswitch {s_{\\scrs\\PW}}^4}\n\\sum_f c_f \\left(1+\\frac{c_{2\\beta}}{s_{2\\beta}}\\xi_\\mathswitchr {A_0}^f\\right) m_f^2\n\\left[\\mathswitch {M_\\PAO}^2-2\\mathswitch {M_\\PHP}^2+\\frac{s_{2\\alpha}}{s_{2\\beta}}(\\mathswitch {M_\\Ph}^2-\\mathswitch {M_\\PH}^2)\\right]\n\\end{align}\nwith the massive fermions $f = \\mathswitchr e, \\dots,\\mathswitchr t$ and the generation index $i$. For the bosonic contribution we obtain\n\\begin{align}\n\\label{eq:L3bos}\n\\delta \\lambda_3\\big|_\\mathrm{UV,bos}\n= \\,& \\Delta_\\mathrm{UV}\\frac{1}{32 \\pi^2} \\big[(\\lambda_1 + \\lambda_2) (6 \\lambda_3 + 2 \\lambda_4) + 4 \\lambda_3^2 + 2 \\lambda_4^2 + 2 \\lambda_5^2\n-3 \\lambda_3 (3 g_2^2 + g_1^2 )\n\\nonumber\\\\\n& +\\textstyle\\frac{3}{4} (3 g_2^4 + g_1^4 - 2 g_2^2 g_1^2)\n\\big]\n\\nonumber\\\\\n= &\n\\,\\Delta_\\mathrm{UV} \\frac{\\lambda_5^2}{2 \\pi^2 s_{2\\beta}^2}\n+ \\Delta_\\mathrm{UV} \\frac{e^2 \\lambda_5}{128 \\mathswitch {M_\\PW}^2 \\pi^2 \\mathswitch {s_{\\scrs\\PW}}^2 s_{2\\beta}^3}\\Big[\n12 (2 \\mathswitch {M_\\PW}^2 + \\mathswitch {M_\\PZ}^2) s_{2 \\beta}^3 \n+ \\mathswitch {M_\\PAO}^2 (27 s_{2 \\beta} - s_{6 \\beta}) \n\\nonumber\\\\\n&+ 2 \\mathswitch {M_\\PHP}^2 (-19 s_{2 \\beta} + s_{6 \\beta}) \n + \\mathswitch {M_\\PH}^2 (-22 s_{2 \\alpha} - 3 s_{2 (\\alpha - 2 \\beta)} - 8 s_{2 \\beta} + s_{2 (\\alpha + 2 \\beta)}) \n\\nonumber\\\\\n& + \\mathswitch {M_\\Ph}^2 (22 s_{2 \\alpha} + 3 s_{2 (\\alpha - 2 \\beta)} - 8 s_{2 \\beta} - s_{2 (\\alpha + 2 \\beta)}) \\Big]\n\\nonumber\\\\\n& + \\Delta_\\mathrm{UV} \\frac{e^4}{256 \\mathswitch {M_\\PW}^4 \\pi^2 \\mathswitch {s_{\\scrs\\PW}}^4 s_{2\\beta}^3} \\Big[\n-2 \\mathswitch {M_\\PH}^4 (-3 + c_{2 (\\alpha - \\beta)} + 2 c_{2 (\\alpha + \\beta)}) s_{2 \\alpha} \n\\nonumber\\\\\n& + 4 \\mathswitch {M_\\Ph}^2 \\mathswitch {M_\\PH}^2 (c_{2 (\\alpha - \\beta)} + 2 c_{2 (\\alpha + \\beta)}) s_{2 \\alpha}\n- 2 \\mathswitch {M_\\Ph}^4 (3 + c_{2 (\\alpha - \\beta)} + 2 c_{2 (\\alpha + \\beta)}) s_{2 \\alpha} \n\\nonumber\\\\\n& - \\mathswitch {M_\\PAO}^4 (-7 s_{2 \\beta} + s_{6 \\beta}) \n- \\mathswitch {M_\\PAO}^2 \\mathswitch {M_\\PH}^2 (11 s_{2 \\alpha} + s_{2 (\\alpha - 2 \\beta)} + 2 s_{2 \\beta})\n\\nonumber\\\\\n& + \\mathswitch {M_\\PAO}^2 \\mathswitch {M_\\Ph}^2 (11 s_{2 \\alpha} + s_{2 (\\alpha - 2 \\beta)} - 2 s_{2 \\beta})\n+ 12 \\mathswitch {M_\\PHP}^4 s_{2 \\beta}^3 \n+ 16 \\mathswitch {M_\\PH}^2 \\mathswitch {M_\\PHP}^2 s_{2 \\beta} s_{\\alpha + \\beta}^2\n\\nonumber\\\\\n&+ 16 \\mathswitch {M_\\Ph}^2 \\mathswitch {M_\\PHP}^2 c_{\\alpha + \\beta}^2 s_{2 \\beta}\n+ 2 \\mathswitch {M_\\PAO}^2 \\mathswitch {M_\\PHP}^2 (-11 s_{2 \\beta} + s_{6 \\beta})\n\\nonumber\\\\\n&+ 6 (2 \\mathswitch {M_\\PW}^2 + \\mathswitch {M_\\PZ}^2) s_{2 \\beta}^2 ((\\mathswitch {M_\\Ph}^2 - \\mathswitch {M_\\PH}^2) s_{2 \\alpha} + (\\mathswitch {M_\\PAO}^2 - 2 \\mathswitch {M_\\PHP}^2) s_{2 \\beta})\n\\nonumber\\\\\n&+ 6 (6 \\mathswitch {M_\\PW}^4 - 4 \\mathswitch {M_\\PW}^2 \\mathswitch {M_\\PZ}^2 + \\mathswitch {M_\\PZ}^4) s_{2 \\beta}^3\n\\Big].\n\\end{align}\n\n\\subsubsection{The FJ tadpole scheme}\n\\label{sec:FJscheme}\n\nSince tadpole loop contribution $T_S$ are gauge dependent \\cite{Degrassi:1992ff}, the connection among bare parameters potentially becomes gauge dependent \nif $\\delta t_S=-T_S$ enters the relations between bare parameters, as it is the case if renormalized tadpole parameters $t_S$ \nare forced to vanish. Note that these gauge dependences systematically cancel if on-shell renormalization conditions are employed, i.e.\\ if predictions \nfor observables are parameterized by directly measurable input parameters. If some input parameters are renormalized in the {$\\overline{\\mathrm{MS}}$}{} scheme this cancellation of gauge dependences does not take place anymore in general, and the gauge dependence is ma\\-ni\\-fest \nin relations between predicted observables and input parameters at NLO.\nIn the $\\overline{\\mathrm{MS}}(\\alpha)$ and the $\\overline{\\mathrm{MS}}(\\lambda_3)$ renormalization schemes, the bare de\\-finitions of $\\alpha$ and $\\beta$ contain tadpole terms \nleading to a gauge dependence (although the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme is gauge \nindependent at NLO in $R_\\xi$ gauges).\n \nFleischer and Jegerlehner \\cite{Fleischer:125234} proposed a renormalization scheme for the SM, referred to as the FJ scheme in the following, that preserves gauge independence for all bare parameters, including the \nmasses and mixing angles%\n\\footnote{A similar scheme, called $\\beta_h$~scheme, was suggested in \\citere{Actis:2006ra}.\nA comparison of that approach to the conventional {$\\overline{\\mathrm{MS}}$}{} and FJ schemes can be found in\n\\citere{Denner:2016etu}.}.\nIn this scheme, the parameters are defined in such a way that tadpole terms do not enter the definition of any bare parameter so that all relations among bare parameters remain gauge independent. \nThis can be achieved by demanding that {\\it bare} tadpole terms vanish,\n$t_{S,0}=0$,\nfor all fields $S$ with the quantum numbers of the vacuum. Since tadpole conditions have no effect on physical observables and change only the bookkeeping, such a procedure is possible. The disadvantage is that now tadpole diagrams have to be taken explicitly into account \nin all higher-order calculations. In particular, the one-particle reducible tadpole contributions destroy the simple relation between propagators and \ntwo-point functions. In the SM, the FJ scheme does not affect observables if all parameters are \nrenormalized using on-shell conditions---as usually done---except \nfor the strong coupling constant $\\alpha_\\mathrm{s}$, which is, however, directly related to the strong gauge coupling, a model defining parameter.\nA gauge-independent renormalization scheme for the THDM can be defined by applying the FJ prescription and imposing the {$\\overline{\\mathrm{MS}}$}{} condition on mixing angles \\cite{Denner:2016etu, Krause:2016oke}.\nThe {\\it bare} physical parameters defined in the FJ scheme \ndiffer by NLO tadpole contributions (including divergent and finite terms)\nfrom the gauge-dependent definition of the bare parameters $\\{p_{\\mathrm{mass}}\\}$ given in Eq.~\\eqref{eq:physparaNLO}. \nExceptions are $e$ and the parameter $\\lambda_5$, which is a parameter of the basic potential and therefore gauge independent by construction. \nThe renormalization of $\\lambda_5$\nin {$\\overline{\\mathrm{MS}}$}{} is identical to the one in the previous schemes.\n\nIt should be noted that in Refs.~\\cite{Denner:2016etu, Krause:2016oke} $m_{12}^2$ is chosen as independent\nparameter in contrast to our choice of $\\lambda_5$. The latter, however, is closer to common practice used in\nthe MSSM~\\cite{Frank2007,Baro:2008bg}.\nMoreover, in Refs.~\\cite{Krause:2016oke,Denner:2016etu}\ntadpole counterterms are reintroduced by shifting the Higgs fields according to $\\eta_i \\to \\eta_i + \\Delta v_{\\eta, i}$, $i = 1,2$, where \nthe constants $\\Delta v_{\\eta, i}$ can be chosen arbitrarily, since physical observables do not depend on this shift,\nwhich can be interpreted as an unobservable change of the integration variables in the path integral. In Refs.~\\cite{Krause:2016oke,Denner:2016etu}, this freedom of choice is exploited, \nand $\\Delta v_{\\eta, i}$ \nare chosen in such a way that the fields $\\eta_1$, $\\eta_2$ do not develop vevs at all order\n. This affects the form of the counterterm Lagrangian and the definition of the renormalization constants with the consequence that the formulae given in Eq.~\\eqref{eq:dVH2p} and Sect.\\ref{sec:onshellrenconditions} cannot be applied.\n\nWe have implemented the FJ scheme following the strategy of Ref.~\\cite{Denner:2016etu} by performing the shifts $\\eta_i \\to \\eta_i+ \\Delta v_{\\eta, i}$ and in an alternative, simpler (but physically equivalent) way. \nIn this simplified approach we keep the dependence of the Lagrangian in terms of gauge-dependent masses and couplings. In addition we keep the tadpole renormalization \ncondition~\\eqref{eq:tadpoleren}, so that the definitions of the renormalization constants of the on-shell parameters and the $Z$ factors according to Sect.~\\ref{sec:onshellrenconditions} remain valid (otherwise we needed to take into account actual tadpole diagrams everywhere).\nIn this simplified approach the counterterms for $\\alpha$ and $\\beta$ which reproduce the\nresults in the FJ scheme result from the previously derived $\\delta\\alpha$ and $\\delta\\beta$\nby adding appropriate finite terms,\n\\begin{align}\n\\delta \\alpha\\big|_\\mathrm{FJ}&= \\delta \\alpha+\\mbox{finite terms},\\nonumber\\\\\n\\delta \\beta\\big|_\\mathrm{FJ} &= \\delta \\beta +\\mbox{finite terms},\n\\end{align}\nwhich depend on the (finite parts of the) tadpole contributions $T_\\mathswitchr H$ and $T_\\mathswitchr h$.\n\nBefore performing the full calculations, we outline the strategy of the derivation of\nthose finite terms for $\\beta$; for $\\alpha$ everything works analogously.\nWe start by exploiting the fact that the form of the tadpole renormalization cannot\nchange physical results if all counterterms for independent parameters are determined \nby the same physical conditions. This means, as mentioned above, that we can simply define\nthe bare tadpoles to vanish, but this forces us to include all explicit tadpole\ncontributions to Green functions. \nWe indicate quantities in this variant by a superscript ``$t$'' in the following. \nWe get the same physical predictions in this ``$t$-variant'' if we use the\ncounterterm\n\\begin{align}\n\\delta \\beta^t &= \\delta \\beta +\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)\n\\label{eq:dbt_def}\n\\end{align}\ninstead of $\\delta\\beta$,\nwhere $\\delta\\beta^t$ is calculated in the same way as $\\delta \\beta$, but with\ntadpole counterterms omitted and \nexplicit tadpole diagrams (including divergent and finite parts)\nin the occurring Green functions taken into account.\nNote that the {$\\overline{\\mathrm{MS}}$}{} prescription to include only divergent terms, which is employed to define\n$\\delta \\beta$, is not applied to the new tadpole contribution $\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)$.\nOtherwise the new $\\Delta\\beta^t$ terms could not be fully compensated by explicit \ntadpole contributions occuring elsewhere, so that\nthere would be differences in the renormalized amplitudes.\nIn fact, applying the {$\\overline{\\mathrm{MS}}$}{} prescription to $\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)$ as well\ndefines the FJ renormalization scheme,\n\\begin{align}\n\\delta \\beta^t\\big|_\\mathrm{FJ} &= \\delta \\beta +\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)\\big|_\\mathrm{UV}.\n\\end{align}\nThe quantity $\\delta \\beta^t\\big|_\\mathrm{FJ}$ is the gauge-independent counterterm for $\\beta$\nintroduced in Ref.~\\cite{Denner:2016etu} which is to be used in the $t$-variant,\nwhere all explicit tadpole diagrams are included in Green functions \n(or equivalently are redistributed by the $\\Delta v$ shift as described Ref.~\\cite{Denner:2016etu}).\nWe can translate the FJ renormalization prescription back to our \nrenormalization scheme (with vanishing renormalized tadpoles) by the counterpart\nof Eq.~\\refeq{eq:dbt_def}, but now formulated in the FJ scheme,\n\\begin{align}\n\\delta \\beta^t\\big|_\\mathrm{FJ} &= \\delta \\beta\\big|_\\mathrm{FJ} +\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h),\n\\label{eq:dbtFJ_def}\n\\end{align}\ni.e.\\ $\\delta \\beta\\big|_\\mathrm{FJ}$ is the counterterm for $\\beta$ to be used in our\ncounterterm Lagrangian in order to calculate renormalized amplitudes in the FJ scheme.\nCombining the above formulas, we obtain the finite difference between\n$\\delta \\beta$ in the (gauge-dependent) {$\\overline{\\mathrm{MS}}$}{} scheme and \n$\\delta \\beta\\big|_\\mathrm{FJ}$ in the (gauge-independent) FJ scheme,\n\\begin{align}\n\\delta \\beta\\big|_\\mathrm{FJ} &= \\delta \\beta^t\\big|_\\mathrm{FJ} -\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)\n\\nonumber\\\\\n&= \\delta \\beta +\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)\\big|_\\mathrm{UV} -\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)\n\\nonumber\\\\\n&= \\delta \\beta -\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)\\big|_\\mathrm{finite}.\n\\end{align}\n\n\\paragraph{The renormalization constant \\boldmath{$\\delta \\beta\\big|_\\mathrm{FJ}$}:}\nWe begin our calculation of $\\delta\\beta|_\\mathrm{FJ}$ with an alternative \ncomputation of $\\delta \\beta$ in the {$\\overline{\\mathrm{MS}}$}{}$(\\alpha)$ scheme, because \nEq.~\\eqref{eq:vacrencond} cannot be applied in the FJ scheme.\nTo avoid the use of Eq.~\\eqref{eq:vacrencond},\nwe calculate the counterterm in the {$\\overline{\\mathrm{MS}}$}{}$(\\alpha)$ scheme from the field \nrenormalization constant by employing the last two equations of Eq.~\\eqref{eq:fieldrenconstrel}. \nThis results in\n\\begin{align}\n \\delta \\beta=\\frac{1}{4} \\left(\\delta Z_{\\mathswitchr {G_0}\\mathswitchr {A_0}}-\\delta Z_{\\mathswitchr {A_0}\\mathswitchr {G_0}}\\right)\\big|_\\mathrm{UV}\n=\\left.\\frac{\n2\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2\n-\\mathrm{Re}\\,\\Sigma^{\\mathswitchr {G_0}\\mathswitchr {A_0}}(\\mathswitch {M_\\PAO}^2)-\\mathrm{Re}\\,\\Sigma^{\\mathswitchr {G_0}\\mathswitchr {A_0}}(0)\n}{2 \\mathswitch {M_\\PAO}^2}\\right|_\\mathrm{UV},\n \\label{eq:betaMSaltern}\n\\end{align}\nwith $\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2$ as given in Eq.~\\eqref{eq:Mabchoice2}. In the second step, $\\delta Z_{\\mathswitchr {G_0}\\mathswitchr {A_0}}$ from Eq.~\\eqref{eq:mixfieldrendef} and\n\\begin{align}\n\\delta Z_{\\mathswitchr {A_0}\\mathswitchr {G_0}}=2 \\, \\frac{-\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2+\\mathrm{Re}\\,\\Sigma^{\\mathswitchr {G_0}\\mathswitchr {A_0}}(0)}{\\mathswitch {M_\\PAO}^2}\n\\label{eq:dZAG}\n\\end{align}\nhave been used. Equation~\\refeq{eq:dZAG} results from demanding finiteness of the\n$G_0A_0$ mixing self-energy at zero-momentum transfer, $k^2=0$, but actually any other \nvalue of $k^2$ would be possible as well, since we only have to remove all UV-divergent terms\nin the mixing. \nThe non-vanishing tadpole counter\\-terms in the {$\\overline{\\mathrm{MS}}$}{}$(\\alpha)$ scheme are $\\delta t_S=-T_S$.\nIn the transition to the $t$-variant, $\\delta \\beta$ gets modified by two kind of terms:\nFirst, there are no tadpole counterterms, i.e.\\ the $\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2$ term is absent,\nand second, there are explicit tadpole contributions to $\\Sigma^{\\mathswitchr {G_0}\\mathswitchr {A_0}}$.\nThis implies\n\\begin{align}\n\\Delta\\beta^t(T_\\mathswitchr H,T_\\mathswitchr h) &= \\delta\\beta^t-\\delta\\beta\n= -\\frac{\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2}{\\mathswitch {M_\\PAO}^2} \n-\\left.\\frac{\\mathrm{Re}\\,\\Sigma^{t,\\mathswitchr {G_0}\\mathswitchr {A_0}}(\\mathswitch {M_\\PAO}^2)+\\mathrm{Re}\\,\\Sigma^{t,\\mathswitchr {G_0}\\mathswitchr {A_0}}(0)}{2 \\mathswitch {M_\\PAO}^2}\\right|_\\mathrm{T_\\mathswitchr H,T_\\mathswitchr h},\n\\end{align}\nwhere the superscript ``$t$'' indicates that one-particle-reducible tadpole diagrams are \nincluded in the self-energies.\nThe subscript ``$T_\\mathswitchr H,T_\\mathswitchr h$'' means that only the those explicit tadpole contributions\nare taken into account here.\nInserting $\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2$ from Eq.~\\refeq{eq:Mabchoice2}\nand evaluating the (momentum-independent) tadpole diagrams for the $G_0A_0$ mixing,\n$\\Delta\\beta^t$ evaluates to\n\\begin{align}\n\\lefteqn{\\Delta\\beta^t(T_\\mathswitchr H,T_\\mathswitchr h)} \n\\nonumber\\\\*\n&= -\\frac{1}{\\mathswitch {M_\\PAO}^2} \\left[\n\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2\n+\\left.\\mathrm{Re}\\,\\Sigma^{t,\\mathswitchr {A_0}\\mathswitchr {G_0}}(0)\\right|_\\mathrm{T_\\mathswitchr H,T_\\mathswitchr h} \\right]\n\\nonumber\\\\\n&=-\\frac{1}{ \\mathswitch {M_\\PAO}^2}\\left[\n\\delta M_{\\mathswitchr {G_0}\\mathswitchr {A_0}}^2\n+ \\parbox[c][1cm]{2.8cm}{\\includegraphics[scale=1]{.\/diagrams\/tadpoleSEHAG1.eps}}+\\parbox[c][1cm]{2.8cm}{\\includegraphics[scale=1]{.\/diagrams\/tadpoleSEhAG2.eps}}\\right]\n\\nonumber\\\\\n&=-\\frac{e}{2 \\mathswitch {M_\\PW} \\mathswitch {s_{\\scrs\\PW}}\\mathswitch {M_\\PAO}^2} \\left[\nT_\\mathswitchr H s_{\\alpha -\\beta }\n+T_\\mathswitchr h c_{\\alpha -\\beta }\n+ T_\\mathswitchr H \\frac{\\left(\\mathswitch {M_\\PAO}^2-\\mathswitch {M_\\PH}^2\\right) s_{\\alpha -\\beta }}{\\mathswitch {M_\\PH}^2} \n+T_\\mathswitchr h \\frac{\\left(\\mathswitch {M_\\PAO}^2-\\mathswitch {M_\\Ph}^2\\right) c_{\\alpha -\\beta }}{\\mathswitch {M_\\Ph}^2}\n\\right]\n\\nonumber\\\\\n&= -\\frac{e}{2 \\mathswitch {s_{\\scrs\\PW}} \\mathswitch {M_\\PW}} \\left(\nT_\\mathswitchr H \\frac{ s_{\\alpha -\\beta }}{\\mathswitch {M_\\PH}^2}\n+T_\\mathswitchr h \\frac{ c_{\\alpha -\\beta }}{\\mathswitch {M_\\Ph}^2} \n\\right).\n\\end{align}\nThe counterterm $\\delta\\beta^t\\big|_\\mathrm{FJ}$ of the FJ scheme in the $t$-variant,\nthus, reads\n\\begin{align}\n\\delta \\beta^t\\big|_\\mathrm{FJ} &= \\delta \\beta +\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)\\big|_\\mathrm{UV}\n= \\delta \\beta \n-\\left.\\frac{e}{2 \\mathswitch {s_{\\scrs\\PW}} \\mathswitch {M_\\PW}} \\left(\nT_\\mathswitchr H \\frac{ s_{\\alpha -\\beta }}{\\mathswitch {M_\\PH}^2}\n+T_\\mathswitchr h \\frac{ c_{\\alpha -\\beta }}{\\mathswitch {M_\\Ph}^2} \n\\right)\\right|_\\mathrm{UV},\n\\label{eq:dbtFJ_expl}\n\\end{align}\nwhich is in agreement with Ref.~\\cite{Krause:2016oke,Denner:2016etu}. \nThis translates to our treatment of tadpoles as\n\\begin{align}\n\\label{eq:dbFJ_expl}\n\\delta \\beta\\big|_\\mathrm{FJ} &= \\delta \\beta \n-\\Delta\\beta^t(T_\\mathswitchr H ,T_\\mathswitchr h)\\big|_\\mathrm{finite}\n= \\delta \\beta\n+ \\left.\\frac{e}{2 \\mathswitch {s_{\\scrs\\PW}} \\mathswitch {M_\\PW}} \\left(\nT_\\mathswitchr H \\frac{ s_{\\alpha -\\beta }}{\\mathswitch {M_\\PH}^2}\n+T_\\mathswitchr h \\frac{ c_{\\alpha -\\beta }}{\\mathswitch {M_\\Ph}^2} \n\\right)\\right|_\\mathrm{finite},\n\\end{align}\nwhere again ``finite'' means that $\\Delta_\\mathrm{UV}$ is set to zero in the tadpole contribution.\nUsing this counterterm, it is possible to keep the form of the\ncounterterm Lagrangian derived \nin Sect.~\\ref{sec:CTLagrangian}\nto obtain results in the gauge-independent\nFJ scheme, although the above counterterm Lagrangian employs a gauge-dependent\n(but very convenient) tadpole renormalization.\n\nAlternatively, $\\delta\\beta$ could be fixed by an analogous consideration\nof the $\\mathswitchr Z\\mathswitchr {A_0}$ mixing, leading to \n\\begin{align}\n\\delta\\beta|_\\mathrm{UV}\n=\\left[ \\frac{s_{2\\beta}}{4c_{2\\alpha}}\\left(\\delta Z_{\\mathswitchr H}-\\delta Z_{\\mathswitchr h}\\right) \n+ \\frac{\\Sigma^{\\mathswitchr Z\\mathswitchr {A_0}}(k^2)}{\\mathswitch {M_\\PZ}} \\right]_\\mathrm{UV},\n\\end{align}\nwhich is independent of $k^2$ and does not use relation~\\refeq{eq:vacrencond}.\nThe transition to the FJ scheme then simply amounts to replacing the\none-particle-irreducible self-energy $\\Sigma^{\\mathswitchr Z\\mathswitchr {A_0}}$ by\n$\\Sigma^{t,\\mathswitchr Z\\mathswitchr {A_0}}$, which includes tadpole diagrams. \nThe result $\\delta \\beta^t\\big|_\\mathrm{FJ}$ of this procedure\nis again given by Eq.~\\refeq{eq:dbtFJ_expl}, as it should be.\n\n\\paragraph{The renormalization constant \\boldmath{$\\delta \\alpha\\big|_\\mathrm{FJ}$}:}\nWe apply the same method to the renormalization constant \n$\\delta \\alpha\\big|_\\mathrm{FJ}$, starting from Eq.~\\refeq{eq:deltaalphaderiv2}.\nThe difference between $\\delta \\alpha$ and $\\delta \\alpha^t$ is entirely given\nby the explicit tadpole diagrams that appear in the change from \n$\\Sigma^{\\mathswitchr H\\mathswitchr h}$ to $\\Sigma^{t,\\mathswitchr H\\mathswitchr h}$ in Eq.~\\refeq{eq:deltaalphaderiv2},\n\\begin{align}\n\\Delta\\alpha^t(T_\\mathswitchr H,T_\\mathswitchr h) &= \\delta\\alpha^t-\\delta\\alpha\n= \\left. \\mathrm{Re}\\frac{\\Sigma^{t,\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\PH}^2)+\\Sigma^{t,\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\Ph}^2)}{2(\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2)}\n\\right|_\\mathrm{T_\\mathswitchr H,T_\\mathswitchr h},\n\\end{align}\nwhich evaluates to\n\\begin{align}\n\\Delta\\alpha^t(T_\\mathswitchr H,T_\\mathswitchr h) \n& = \\left. \\mathrm{Re}\\frac{\\Sigma^{t,\\mathswitchr H\\mathswitchr h}(\\mathswitch {M_\\Ph}^2)}{\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2} \\right|_\\mathrm{T_\\mathswitchr H,T_\\mathswitchr h}\n=\\frac{1}{\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2}\\left[\\parbox[c][1cm]{2.5cm}{\\includegraphics[scale=1]{.\/diagrams\/tadpoleSEhHh2.eps}}+\n\\parbox[c][1cm]{2.5cm}{\\includegraphics[scale=1]{.\/diagrams\/tadpoleSEHHh1.eps}} \\right]\n\\nonumber\\\\\n&=\\frac{e}{\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2} \\left( \nT_\\mathswitchr H \\frac{C_{\\mathswitchr h\\mathswitchr H\\PH}}{\\mathswitch {M_\\PH}^2}\n+T_\\mathswitchr h \\frac{C_{\\mathswitchr h\\Ph\\mathswitchr H}}{\\mathswitch {M_\\Ph}^2} \n\\right),\n\\end{align}\nwith the coupling factors of the $\\mathswitchr h\\mathswitchr H\\PH$ and $\\mathswitchr h\\Ph\\mathswitchr H$ vertices\n\\begin{subequations}\n\\begin{align}\nC_{\\mathswitchr h\\mathswitchr H\\PH} &= \\frac{e s_{\\beta-\\alpha} }{2\\mathswitch {M_\\PW}\\mathswitch {s_{\\scrs\\PW}} s_{2\\beta}}\\biggl[ \n-(3s_{2\\alpha}+s_{2\\beta})\n\\left(\\mathswitch {M_\\PAO}^2 +4 \\lambda_5 \\frac{\\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2}{e^2}\\right)\n+ s_{2 \\alpha} \\left(\\mathswitch {M_\\Ph}^2 + 2 \\mathswitch {M_\\PH}^2\\right) \\biggr], \n\\\\\nC_{\\mathswitchr h\\Ph\\mathswitchr H} &= \\frac{ec_{\\beta-\\alpha}}{2\\mathswitch {M_\\PW}\\mathswitch {s_{\\scrs\\PW}} s_{2\\beta}} \\biggl[\n(3s_{2\\alpha}-s_{2\\beta})\n\\left( \\mathswitch {M_\\PAO}^2 + 4 \\lambda_5 \\frac{\\mathswitch {M_\\PW}^2 \\mathswitch {s_{\\scrs\\PW}}^2}{e^2}\\right)\n- s_{2\\alpha}\\left(2 \\mathswitch {M_\\Ph}^2 + \\mathswitch {M_\\PH}^2 \\right) \\biggr].\n\\end{align}\n\\end{subequations}\nThe counterterm $\\delta\\alpha^t\\big|_\\mathrm{FJ}$ of the FJ scheme in the $t$-variant,\nthus, reads\n\\begin{align}\n\\delta \\alpha^t\\big|_\\mathrm{FJ} &= \\delta \\alpha +\\Delta\\alpha^t(T_\\mathswitchr H ,T_\\mathswitchr h)\\big|_\\mathrm{UV}\n= \\delta \\alpha \n+\\frac{e}{\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2} \\left( \nT_\\mathswitchr H \\frac{C_{\\mathswitchr h\\mathswitchr H\\PH}}{\\mathswitch {M_\\PH}^2}+\nT_\\mathswitchr h \\frac{C_{\\mathswitchr h\\Ph\\mathswitchr H}}{\\mathswitch {M_\\Ph}^2} \n\\right)\\bigg|_\\mathrm{UV},\n\\label{eq:datFJ_expl}\n\\end{align}\nwhich is again in agreement with Ref.~\\cite{Krause:2016oke,Denner:2016etu}. \nThis translates to our treatment of tadpoles as\n\\begin{align}\n\\delta \\alpha\\big|_\\mathrm{FJ} &= \\delta \\alpha \n-\\Delta\\alpha^t(T_\\mathswitchr H ,T_\\mathswitchr h)\\big|_\\mathrm{finite}\n= \\delta \\alpha\n+\\frac{e}{\\mathswitch {M_\\Ph}^2-\\mathswitch {M_\\PH}^2} \\left( \nT_\\mathswitchr H \\frac{C_{\\mathswitchr h\\mathswitchr H\\PH}}{\\mathswitch {M_\\PH}^2}\n+T_\\mathswitchr h \\frac{C_{\\mathswitchr h\\Ph\\mathswitchr H}}{\\mathswitch {M_\\Ph}^2} \n\\right)\\bigg|_\\mathrm{finite}.\n\\end{align}\nConcerning the use of $\\delta \\alpha\\big|_\\mathrm{FJ}$ in our counterterm Lagrangian\nto obtain renormalized amplitudes in the gauge-independent FJ scheme,\nthe same comments made above for $\\delta\\beta\\big|_\\mathrm{FJ}$ apply.\n\n\\subsubsection{\\boldmath{The FJ($\\lambda_3$) scheme}}\n\\label{sec:FJL3scheme}\n\nIn the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme, the parameters $\\lambda_{3}$ \nand $\\lambda_{5}$ are defining parameters of the basic para\\-meterization and gauge independent by construction. Therefore, the condition on $\\delta \\beta$ is the only renormalization condition potentially being gauge dependent. To provide a fully\ngauge-independent renormalization scheme where $\\lambda_3$ is an independent quantity, we apply the FJ scheme to the parameter $\\beta$ and keep the renormalization of \n$\\lambda_{3}$ and $\\lambda_{5}$ \nas in the $\\overline{\\mathrm{MS}}(\\lambda_3)$. We call the resulting scheme the FJ($\\lambda_3$) scheme. The renormalization of the parameters reads:\n\\begin{align}\n \\delta \\beta\\big|_\\mathrm{FJ}\\text{ as in Eq.~\\eqref{eq:dbFJ_expl}},\\nonumber\\\\\n \\delta \\lambda_3 \\text{ as in Eqs.~\\eqref{eq:deltaL3}}{-}\\refeq{eq:L3bos}.\\nonumber\n\\end{align}\n\n\\subsection{Conversion between different renormalization schemes}\n\\label{sec:inputconversion}\n\nIn the previous section, we have\npresented four different renormalization schemes, which treat the mixing parameters differently. When observables calculated in different renormalization schemes are compared, particular care has to be taken that the input parameters are consistently translated from one scheme to the other. The bare values of identical independent\nparameters are equal and independent of the renormalization scheme. \nExemplarily, for a parameter $p$, the renormalized values $p^{(1)}$ and $p^{(2)}$ in two different renormalization schemes~1 and 2 are connected via the bare parameter $p_0$, \n\\begin{align}\np_0=p^{(1)}+\\delta p^{(1)} (p^{(1)})=p^{(2)} +\\delta p^{(2)} (p^{(2)}),\n\\end{align}\nwithin the considered order. \nIf $p$ is a dependent parameter in one or both schemes, it must be calculated from the independent renormalized para\\-meters and their counterterms from the relations between bare and renormalized quantities. For converting an input value from one scheme to another, one can solve for one renormalized quantity\n\\begin{align}\n\\label{eq:convertinputparam}\n p^{(1)}=p^{(2)} +\\delta p^{(2)} (p^{(2)})-\\delta p^{(1)} (p^{(1)}).\n\\end{align}\nAt NLO, this equation can be linearized by substituting the input value of $p^{(1)}$ by $p^{(2)}$ \nin the computation of the last counterterm. The differences to an exact solution are of higher order and beyond our desired NLO\naccuracy. However, large counterterms or small tree-level values can spoil the approximation so that in this case a proper solution using numerical techniques could improve the results. \nAnother benefit of a full solution of the implicit equation is the possibility that one can \nswitch to another scheme and back in a self-consistent way, while start and end scenarios\nin scheme~(1) do not exactly coincide when switching from scheme~(1) to (2) and back to (1)\nusing the linearized approximation.\nThe comparison of both methods allows for a consistency check of the computation and for an analysis of perturbative stability.\nWe have derived the Higgs mixing angles $\\alpha$ and $\\beta$ and their counterterm in all schemes. The finite parts of the gauge-dependent counterterms $\\delta \\alpha$, $\\delta \\beta$ are given here for the different renormalization schemes indicated by the respective index: \n\\begin{subequations}\n\\label{eq:fintermsalpha}\n\\begin{align}\n \\delta \\alpha\\big|_{\\overline{\\mathrm{MS}}(\\alpha),\\mathrm{finite}} &= 0,\\\\\n \\delta \\alpha\\big|_{\\overline{\\mathrm{MS}}(\\lambda_3),\\mathrm{finite}} \n&= \\delta \\alpha\\big|_{\\mathrm{FJ}(\\lambda_3),\\mathrm{finite}} \n= \\frac{f_\\alpha\\{\\delta p'_\\mathrm{mass}\\}}{\\mathswitch {M_\\PH}^2-\\mathswitch {M_\\Ph}^2}\\Big|_\\mathrm{finite},\\\\\n \\delta \\alpha\\big|_{\\mathrm{FJ}(\\alpha),\\mathrm{finite}} &= -\\Delta\\alpha^t(T_\\mathswitchr H,T_\\mathswitchr h)\\big|_\\mathrm{finite}\n=\\left.\\frac{e}{\\mathswitch {M_\\Ph}^2-\\mathswitch {M_\\PH}^2} \\left( \nT_\\mathswitchr H \\frac{C_{\\mathswitchr h\\mathswitchr H\\PH}}{\\mathswitch {M_\\PH}^2}\n+T_\\mathswitchr h \\frac{C_{\\mathswitchr h\\Ph\\mathswitchr H}}{\\mathswitch {M_\\Ph}^2} \n\\right)\\right|_\\mathrm{finite}.\n\\end{align}\n\\end{subequations}\nFor the angle $\\beta$ we obtain the following finite terms \nin the $\\overline{\\mathrm{MS}}$ and the FJ schemes,\n\\begin{subequations}\n\\label{eq:fintermsbeta}\n\\begin{align}\n \\delta \\beta\\big|_{\\overline{\\mathrm{MS}(\\alpha)},\\mathrm{finite}} \n&= \\delta \\beta\\big|_{\\overline{\\mathrm{MS}}(\\lambda_3),\\mathrm{finite}} \n= 0,\\\\\n\\delta \\beta\\big|_{\\mathrm{FJ}(\\alpha),\\mathrm{finite}} \n&= \\delta \\beta\\big|_{\\mathrm{FJ(\\lambda_3)},\\mathrm{finite}} \n = -\\Delta\\beta^t(T_\\mathswitchr H,T_\\mathswitchr h)\\big|_\\mathrm{finite}=\\left.\\frac{e}{2 \\mathswitch {s_{\\scrs\\PW}} \\mathswitch {M_\\PW}} \\left(\nT_\\mathswitchr H \\frac{ s_{\\alpha -\\beta }}{\\mathswitch {M_\\PH}^2}\n+T_\\mathswitchr h \\frac{ c_{\\alpha -\\beta }}{\\mathswitch {M_\\Ph}^2} \n\\right)\\right|_\\mathrm{finite}.\n\\end{align}\n\\end{subequations}\nWith these formulae we can convert\nthe input variables for $\\alpha$ and $\\beta$\neasily into each other. \nFor instance, the conversion of the input values of $\\alpha$ and $\\beta$ \ndefined in the $\\overline{\\mathrm{MS}}(\\alpha)$ scheme into \nthe other renormalization schemes reads\n\\begin{subequations}\n\\begin{align}\n\\alpha\\big|_{\\overline{\\mathrm{MS}}(\\lambda_3)} \n& = \\alpha\\big|_{\\overline{\\mathrm{MS}}(\\alpha)}\n+\\left.\\frac{f_\\alpha\\{\\delta p'_\\mathrm{mass}\\}}{\\mathswitch {M_\\Ph}^2-\\mathswitch {M_\\PH}^2}\\right|_\\mathrm{finite}, &\n\\beta\\big|_{\\overline{\\mathrm{MS}}(\\lambda_3)} &=\n\\beta\\big|_{\\overline{\\mathrm{MS}}(\\alpha)},\n\\label{eq:convertinput1}%\n\\\\[.3em]\n\\alpha\\big|_{\\mathrm{FJ}(\\alpha)}\n& = \\alpha\\big|_{\\overline{\\mathrm{MS}}(\\alpha)}\n+\\Delta\\alpha^t(T_\\mathswitchr H,T_\\mathswitchr h)\\big|_\\mathrm{finite}, &\n\\beta\\big|_{\\mathrm{FJ}(\\alpha)}\n& = \\beta\\big|_{\\overline{\\mathrm{MS}}(\\alpha)}\n+\\Delta\\beta^t(T_\\mathswitchr H,T_\\mathswitchr h)\\big|_\\mathrm{finite},\n\\label{eq:convertinput2}%\n\\\\[.3em]\n\\alpha\\big|_{\\mathrm{FJ}(\\lambda_3)}\n& = \\alpha\\big|_{\\overline{\\mathrm{MS}}(\\alpha)}\n+\\left.\\frac{f_\\alpha\\{\\delta p'_\\mathrm{mass}\\}}{\\mathswitch {M_\\Ph}^2-\\mathswitch {M_\\PH}^2}\\right|_\\mathrm{finite}, & \n\\beta\\big|_{\\mathrm{FJ}(\\lambda_3)}\n& = \\beta\\big|_{\\overline{\\mathrm{MS}}(\\alpha)}\n+\\Delta\\beta^t(T_\\mathswitchr H,T_\\mathswitchr h)\\big|_\\mathrm{finite}.\n\\label{eq:convertinput3}%\n\\end{align}\n\\end{subequations}%\nWithin a given scheme, $\\lambda_3$ and $\\alpha$ can be translated into each other\nusing the tree-level relation~\\refeq{eq:lambda3tree}.\nNote that, thus, the numerical values of $\\alpha$, $\\beta$, and $\\lambda_3$ corresponding\nto a given physical scenario of the THDM are different in different renormalization schemes.\nIn turn, fixing the input values in the four renormalization schemes to the same values\ncorresponds to different physical scenarios.\nIn particular, this means that the ``alignment limit'', in which $s_{\\beta-\\alpha}\\to1$\nso that $\\mathswitchr h$ is SM like \n{(see, e.g., \\citeres{Haber:1995be,Gunion:2002zf,Bernon:2015qea}),} is a notion that depends on the renormalization scheme\n(actually even on the scale choice in a given scheme).%\n\\footnote{In this brief account of results, we do not consider the (phenomenologically\ndisfavoured, though not excluded) possibility that the heavier CP-even Higgs boson $\\mathswitchr H$ is SM-like,\nwhich is discussed in detail in \\citere{Bernon:2015wef}.}\n\nExemplarily, the \nconversions of $c_{\\beta-\\alpha}$ from the {$\\overline{\\mathrm{MS}}$}{}$(\\alpha)$ scheme into the $\\overline{\\mathrm{MS}}(\\lambda_3)$ (green), FJ$(\\alpha)$ (pink), and \nFJ$(\\lambda_3)$ (turquoise) schemes are shown in Fig.~\\ref{fig:plot_Umrechnung-von-alpha-LM}. The results of the transformations \nin the inverse directions are displayed in Fig.~\\ref{fig:plot_Umrechnung-nach-alpha-LM}, and all other conversions can be seen as a combination of the presented ones. \n\\begin{figure}\n\\centering\n\\subfigure[]{\n\\label{fig:plot_Umrechnung-von-alpha-LM}\n\\includegraphics{.\/diagrams\/plot_Umrechnung-von-alpha-LM.eps}\n}\n\\hspace{15pt}\n\\subfigure[]{\n\\label{fig:plot_Umrechnung-nach-alpha-LM}\n\\includegraphics{.\/diagrams\/plot_Umrechnung-nach-alpha-LM.eps}\n}\n\\vspace*{-.5em}\n\\caption{(a) Conversion of the value of $c_{\\beta-\\alpha}$ from $\\overline{\\mathrm{MS}}(\\alpha)$ \nto the $\\overline{\\mathrm{MS}}(\\lambda_3)$ (green), FJ$(\\alpha)$ (pink), and FJ($\\lambda_3$) \nschemes (turquoise) for scenario~A. Panel~(b) shows the conversion to the $\\overline{\\mathrm{MS}}(\\alpha)$ scheme using the same colour coding. \nThe solid lines are obtained \nby solving the implicit equations \\refeq{eq:convertinputparam} numerically, the dashed lines\ncorrespond to the linearized approximation.\nThe phenomenologically relevant region is highlighted in the centre.}\n\\label{fig:plotconversionA}%\n\\end{figure}%\nThe input values (defined before the conversion) correspond to the low-mass scenario called ``A''\nof a THDM of Type~I (based on a benchmark scenario of Ref.~\\cite{Haber2015})\nwith \n\\begin{align}\n\\mathswitch {M_\\Ph}=125 \\text{ GeV}, \\quad\n \\mathswitch {M_\\PH}=300 \\text{ GeV}, \\quad \\mathswitch {M_\\PAO}=\\mathswitch {M_\\PHP}=460 \\text{ GeV}, \\quad \\lambda_5=-1.9,\\quad \\tan \\beta=2. \n\\label{eq:scenario}\n \\end{align}\nSpecifically, scenario~A is a scan in $c_{\\beta-\\alpha}$ in the mass parameterization, \nAa and Ab are points of the scan region used to analyze the scale dependence:\n\\begin{subequations}\n\\begin{align}\n\\mbox{A:} \\quad \\cos{(\\beta-\\alpha)} &= -0.2\\ldots0.2, \\\\\n\\mbox{Aa:} \\quad \\cos{(\\beta-\\alpha)} &= +0.1, \\\\\n\\mbox{Ab:} \\quad \\cos{(\\beta-\\alpha)} &= -0.1. \n\\end{align}\n\\label{eq:cba_A}\n\\end{subequations}\nThe {$\\overline{\\mathrm{MS}}$}{} parameters are defined at the scale\n\\begin{align}\n\\label{eq:centralscale}\n \\mu_0=\\frac{1}{5}(\\mathswitch {M_\\Ph}+\\mathswitch {M_\\PH}+\\mathswitch {M_\\PAO}+2\\mathswitch {M_\\PHP}).\n\\end{align}\nThe motivation for this choice will become clear below.\nThe remaining input parameters for the SM part are given in \nApp.~\\ref{app:smparams}.\nIn both plots, we highlight the phenomenologically relevant region in the centre. \nThe solid lines are the result obtained \nby solving the implicit equations \\refeq{eq:convertinputparam} numerically,\nthe dashed lines correspond to a linearized conversion.\nAll curves show only minor conversion effects in the parameter values, i.e.\\ the \nsolution of the implicit equations\nagrees well with the approximate linearized conversion, affirming that the contributions of the\nhigher-order functions $\\Delta\\alpha^t$, $\\Delta\\beta^t$, and $f_\\alpha$ of \nEqs.~\\refeq{eq:convertinput1}--\\refeq{eq:convertinput3}\nare small, and perturbation theory is applicable. Since the values of the parameters change when going from one renormalization scheme to another, \nthe alignment limit does not persist in these transformations, i.e.\\\nin this scenario the alignment limit sensitively depends on the definition of the parameters at NLO.\n\n{For the schemes with $\\lambda_3$ as input parameter, some singular behaviour in the parameter\nconversion can be observed in the phenomenologically disfavoured region where $c_{\\beta-\\alpha}\\lsim-0.3$.\nThis artifact in the conversion appears when $c_{2\\alpha}\\to0$ (see, e.g., Eq.~\\refeq{eq:MHhdependence}),\nindicating the breakdown of the $\\overline{\\mathrm{MS}}(\\lambda_3)$ and FJ($\\lambda_3$) schemes in such\nparameter regions. Already this case-specific study shows that stability issues of different\nrenormalization schemes have to be carefully carried out for all interesting parameter regions\nand that the applicability of a specific scheme in general does not cover the full THDM \nparameter space, a fact that was also pointed out in \\citere{Krause:2016oke}\nfor the THDM and that is known from NLO calculations in the MSSM (see, e.g., \\citeres{Heinemeyer:2010mm,Chatterjee:2011wc}).\nSpecifically, if the $\\overline{\\mathrm{MS}}(\\lambda_3)$ and FJ($\\lambda_3$) schemes are not\napplicable in a region that might be favoured by future data analyses, it would\nbe desirable and straightforward to replace $\\lambda_3$ by $\\lambda_1$ or $\\lambda_2$\nas independent parameter, thereby defining analogous schemes like $\\overline{\\mathrm{MS}}(\\lambda_1)$, etc..\n\nTo address this issue properly, was our basic motivation to introduce and compare different renormaliztion\nschemes. We will continue this discussion in more detail in a forthcoming publication,\nwhere further THDM scenarios are considered.}\n\n\\section{\\boldmath{The running of the {$\\overline{\\mathrm{MS}}$}{} parameters}}\n\\label{sec:running}\n\nParameters renormalized in the $\\overline{\\mathrm{MS}}$ scheme depend on an unphysical renormalization scale $\\mu_\\mathrm{r}$. The one-loop $\\beta$-function of a parameter $p$ can be obtained from the UV-divergent parts of its \ncounterterm $\\delta p$, \n\\begin{align}\n\\beta_p(\\mu^2_\\mathrm{r})=\\frac{\\partial}{\\partial \\ln{\\mu_\\mathrm{r}^2}} p(\\mu_\\mathrm{r}^2)\n=\\frac{\\partial}{\\partial \\Delta_\\mathrm{UV}} \\delta p.\n\\label{eq:betafct}\n\\end{align}\nSince the renormalization constants are computed in a perturbative manner, the $\\beta$-functions have a perturbative expansion in the coupling parameters.\nNote that the last equality in Eq.~\\refeq{eq:betafct} holds in the FJ schemes\nfor $\\alpha$ and $\\beta$\nonly in the $t$-variant explained above, because the finite contributions\n$\\Delta\\alpha^t$ and $\\Delta\\beta^t$ depend on the scale $\\mu_\\mathrm{r}$.\n\nAs discussed in the previous sections, \nthe ratio of the vevs, $\\tan \\beta$, the Higgs mixing parameter $\\alpha$ or $\\lambda_3$, and the Higgs self-coupling $\\lambda_5$ are renormalized in the $\\overline{\\mathrm{MS}}$ scheme. For each renormalization scheme described in Sect.~\\ref{sec:diffrenschemes}, one obtains a set of coupled RGEs involving the $\\beta$-functions of the independent parameters. Therefore, the scale dependence varies when different schemes are applied. In the perturbative expansion of the $\\beta$-function we consider only the one-loop term, being second order in the coupling constants,\ne.g., in the {$\\overline{\\mathrm{MS}}$}$(\\alpha)$ scheme\n\\begin{align}\n\\beta_p(\\mu^2_\\mathrm{r})= A_p \\alpha_\\mathrm{em}+ B_p \\lambda_5+C_p \\lambda_5^2\/\\alpha_\\mathrm{em}.\n\\end{align}\nThe dependence on the strong coupling constant vanishes at one-loop order as the parameters renormalized in {$\\overline{\\mathrm{MS}}$}{} appear only in couplings of particles that do not interact strongly. The coefficients $A_p, B_p, C_p$ of the respective renormalized parameter can be easily read \nfrom the divergent terms which have been derived in the previous section. We have checked them against the $\\beta$-functions given for $\\lambda_{3}$ and $\\lambda_{5}$ in Ref.~\\cite{Branco:2011iw} and for $\\beta$ in Ref.~\\cite{Sperling:2013eva} (supersymmetric contributions need to be omitted).\n\nIn general, RGEs, which are a set of coupled differential equations, cannot be solved analytically. \nUsually numerical techniques, such as a Runge--Kutta method, need to be employed to solve the RGEs and to compute the values of the parameters at a desired scale.\nMoreover, we emphasize that the renormalization-group flow of a running parameter\ndepends on the renormalization scheme of the full set of independent parameters.\nThat means the fact that we use on-shell quantities, such as all the Higgs-boson masses,\nto fix most of the scalar self-couplings has a significant impact on the running\nof our {$\\overline{\\mathrm{MS}}$}{} parameters.\nThe renormalization-group flow in other schemes was, e.g., investigated in\n\\citeres{Cvetic:1997zd,Branco:2011iw,Bijnens:2011gd,Chakrabarty:2014aya,Ferreira:2015rha}.\n\nThe scale dependence of $c_{\\beta-\\alpha}$ \nfor $\\mu = 100{-}900$~GeV is plotted in Fig.~\\ref{fig:plotrunningA}, for the scenario defined in Eq.~\\eqref{eq:scenario} with $c_{\\beta-\\alpha}=0.1$ (l.h.s) and $c_{\\beta-\\alpha}=-0.1$ (r.h.s) and input values given at the central scale $\\mu_0$ stated in\nEq.~\\refeq{eq:centralscale}. \nWe observe that the choice of the renormalization scheme has a large impact on the scale dependence. While the $\\overline{\\mathrm{MS}}(\\alpha)$ scheme introduces only a mild running, the other schemes show a much stronger scale dependence, so that excluded and unphysical values of input parameters can be reached quickly. A similar observation has also been made in supersymmetric models for the parameter~$\\tan\\beta$~\\cite{Freitas:2002um}. \nGauge-dependent $\\overline{\\mathrm{MS}}$ schemes have a small scale dependence while replacing the parameters by gauge-independent ones like in the FJ schemes introduce additional terms in the $\\beta$-functions which induce a stronger scale dependence.\nIn Fig~\\ref{fig:running_cba-0.1-LM} one can also see that the curves for the $\\overline{\\mathrm{MS}}(\\lambda_3)$ and the FJ($\\lambda_3$) schemes terminate around 250 GeV. \n\\begin{figure}\n \\centering\n \\subfigure[]{\n\\label{fig:running_cba0.1-LM}\n\\includegraphics{.\/diagrams\/running_cba0.1-LM.eps}\n}\n\\hspace{15pt}\n\\subfigure[]{\n\\label{fig:running_cba-0.1-LM}\n\\includegraphics{.\/diagrams\/running_cba-0.1-LM.eps}\n}\n\\vspace*{-.5em}\n \\caption{The running of $c_{\\beta-\\alpha}$ for the low-mass scenario~A with $c_{\\beta-\\alpha}=0.1$ (a) and $c_{\\beta-\\alpha}=-0.1$ (b) in the $\\overline{\\mathrm{MS}}(\\alpha)$ (blue), $\\overline{\\mathrm{MS}}(\\lambda_3)$ (green), FJ($\\alpha)$ (pink), and FJ($\\lambda_3$) (turquoise) schemes.}\n\\label{fig:plotrunningA}\n\\end{figure}\nAt this scale, the running of $\\lambda_3$ yields unphysical values for which \nEq.~\\refeq{eq:lambda3tree}\nwith the given Higgs masses becomes \noverconstrained, and no solution with $|s_{2\\alpha}|\\le1$ exists. This is unique to the $\\lambda_3$ running as only there an \nimplicit equation needs to be solved to obtain the input parameter $\\alpha$. For the other cases we prevent the angles from running out of their domain of definition by solving the running for the tangent function of the angles.\n\n\\section{Implementation into a \\textsc{FeynArts} Model File}\n\\label{sec:FAmodelfile}\nThe \\textsc{Mathematica} package \\textsc{FeynRules} (FR)~\\cite{Christensen2009} \nis a tool to generate Feynman rules from a given \nLagrangian, providing the possibility to produce model files in various output formats\nwhich can be employed by automated amplitude generators. \nWe have inserted the Lagrangian into FR in its internal notation to obtain\nthe corresponding counter\\-term Lagrangian after the renormalization transformations. \nBefore this insertion, we have computed and simplified \nthe Higgs potential and the corresponding counterterm potential~\\eqref{eq:dVH2p} \nwith inhouse \\textsc{Mathematica} routines.\nUsing FR,\nthe tree-level and the counterterm Feynman rules as well as the renormalization conditions \nin the {$\\overline{\\mathrm{MS}}$}{}$(\\alpha)$ and {$\\overline{\\mathrm{MS}}$}{}($\\lambda_3$) schemes have been implemented into \na model file for the amplitude generator \\textsc{FeynArts} (FA)~\\cite{Hahn2001}. \nThe renormalization conditions of the FJ($\\alpha)$ and the FJ($\\lambda_3$) have not \nbeen included in the model file, because using \nEqs.~\\eqref{eq:fintermsalpha},~\\eqref{eq:fintermsbeta} it is straightforward to \nimplement the corresponding finite terms of $\\delta \\alpha$ and $\\delta \\beta$. \nWith such a model file, NLO amplitudes for any process can be generated in an automated way. \n\nThe FA NLO model file for the THDM, obtained with FR, has the following features:\n\\begin{itemize}\n\\item Type I, II, flipped, or lepton-specific THDM;\n\\item all tree-level and counterterm Feynman rules;\n\\item renormalization conditions according to the {$\\overline{\\mathrm{MS}}$}{}$(\\alpha)$ and {$\\overline{\\mathrm{MS}}$}{}$(\\lambda_3)$ schemes;\n\\item all renormalization constants are implemented additionally in $\\overline{\\mathrm{MS}}$ \nas well, which allows for fast checks of UV-finiteness;\n\\item BHS and HK conventions;\n\\item\nCKM matrix set to the unit matrix (the generalization is straightforward).\n\\end{itemize}\nThis model file has been tested intensively, including checks of UV-finiteness for \nseveral processes, both numerically and analytically. \nThis allows for the generation of amplitudes \n(and further processing with \\textsc{FormCalc} \\cite{Hahn1999}) for any process at the one-loop level, at any parameter point of the THDM.\nThe model file can be obtained from the authors upon request.\n\n\\section{\\boldmath{Numerical results for \\boldmath$ \\mathswitchr h \\to \\mathswitchr W\\PW\/\\mathswitchr Z\\PZ \\to 4f$}}\n\\label{se:numerics}\n\nIn this section we present first results from the computation of the decay of the light, neutral \nCP-even Higgs boson of the THDM into four fermions at NLO. The computer program \n\\textsc{Prophecy4f}~\\cite{Bredenstein:2006rh,Bredenstein:2006nk,Bredenstein:2006ha}%\n\\footnote{\\tt http:\/\/prophecy4f.hepforge.org\/index.html}\nprovides a ``\\textbf{PROP}er description of the \\textbf{H}iggs d\\textbf{EC}a\\textbf{Y} into \\textbf{4 F}ermions'' and calculates \nobservables for the decay process $\\mathswitchr h {\\to} \\mathswitchr W\\PW\/\\mathswitchr Z\\PZ {\\to} 4 f$ at NLO EW+QCD in the SM. \nWe have extended this program to the calculation of the corresponding decay in the THDM in such a way that the usage of the program and its applicability as event\ngenerator basically remains the same. \n{Owing to the fact that LO and real-emission amplitudes in the THDM receive only the\nmultiplicative factor $s_{\\beta-\\alpha}$ with respect to the SM, the bremsstrahlung corrections as well\nas the treatment of infrared singularities could be taken over from the SM \ncalculation~\\cite{Bredenstein:2006rh,Bredenstein:2006ha}\nvia simple rescaling.}\nThe calculation in the THDM, the implementation in \\textsc{Prophecy4f}, as well as results of the application will be described in detail in an upcoming publication.\nWe just mention that we employ the complex-mass scheme~\\cite{Denner:2005fg} to describe the \nW\/Z~resonances, as already done in \\citeres{Bredenstein:2006rh,Bredenstein:2006nk,Bredenstein:2006ha}\nfor the SM. Note that the W\/Z-boson masses as well as the weak mixing angle are consistently taken as complex quantities in the complex-mass scheme to guarantee gauge invariance of all amplitudes\nin resonant and non-resonant phase-space regions. \nConsequently all our renormalization constants of the THDM inherit imaginary parts from the complex\ninput values, but the impact of these spurious imaginary parts is beyond NLO and negligible\n(as in the SM). Moreover, we mention that the modified version of \\textsc{Prophecy4f}\nmakes use of the public \\textsc{Collier} library~\\cite{Denner:2016kdg} for the calculation \nof the one-loop integrals.\n{Apart from performing two independent loop calculations, we have verified\nour one-loop matrix elements by numerically comparing our results to\nthe ones obtained in \\citere{Denner:2016etu} for the related\n$\\mathswitchr W\\mathswitchr h\/\\mathswitchr Z\\mathswitchr h$ production channels (including $\\mathswitchr W\/\\mathswitchr Z$ decays)\nusing crossing symmetry.}\n\nIn this paper, we present first results in order to demonstrate the use and the \nself-consistency of our renormalization schemes, employing again the scenario\ninspired by the first benchmark scenario of Ref.~\\cite{Haber2015} where the additional \nHiggs bosons are not very heavy. \nThe input values of the THDM parameters \n{for a Type~I THDM}\nare given in Eqs.~\\refeq{eq:scenario} and \n\\refeq{eq:cba_A}.\nSince $c_{\\beta-\\alpha}$ is the only parameter of the THDM appearing at LO, \nour process is most sensitive to this parameter. We vary $c_{\\beta-\\alpha}$ \nin the range $[-0.2,+0.2]$ in scenario~A for the computation of \nthe partial decay width for $\\mathswitchr h\\to\\mathswitchr W\\PW\/\\mathswitchr Z\\PZ\\to4f$, $\\Gamma^{\\mathswitchr h\\to4f}_\\mathrm{THDM}$,\nwhich is obtained by summing the partial widths of the h~boson\nover all massless four-fermion final states $4f$.\nThe parameters of the SM part of the THDM are collected in App.~\\ref{app:smparams}.\n{Note that a non-trivial CKM matrix would not change our results,\nsince quark mass effects of the first two generations as well as mixing\nwith the third generation are completely negligible in the considered\ndecays.}\n\nTo perform scale variations we take two distinguished points named Aa and Ab with $c_{\\beta-\\alpha}=\\pm0.1$. \nFor the central renormalization scale we use the average mass \n$\\mu_0$ defined in Eq.~\\refeq{eq:centralscale}\nof all scalar degrees of freedom.\nThe scale $\\mu$ of $\\alpha_\\mathrm{s}$ is kept fixed at $\\mu=\\mathswitch {M_\\PZ}$ which is the appropriate scale for the QCD corrections (which are dominated by the hadronic W\/Z decays).\n\n\\subsection{Scale variation of the width}\n\\label{sec:LowMassscalevariation}\n\n{The running of the {$\\overline{\\mathrm{MS}}$}-renormalized parameters $\\alpha$ and $\\beta$ is\ninduced by the Higgs-boson self-energies (and some scalar vertex for $\\lambda_5$),\ni.e.\\ the relevant particles in the loops are all Higgs bosons, the W\/Z bosons,\nand the top quark. If all Higgs-boson masses are near the electroweak scale,\nsay $\\sim 100{-}200\\unskip\\,\\mathrm{GeV}$, \nwhere the W\/Z-boson and top-quark masses are located, \nthen the scale $\\mathswitch {M_\\Ph}$ turns out to be a reasonable scale, as expected.\n{However, if some heavy Higgs-boson masses increase to some generic mass\nscale $M_S$ and the mixing angle\n$\\beta-\\alpha$ stays away from the alignment limit, there is no \ndecoupling of heavy Higgs-boson effects, so that $M_S$ acts\nas generic UV cutoff scale appearing in logarithms $\\log(M_S\/\\mu_\\mathrm{r})$.\nThe renormalization scale $\\mu_\\mathrm{r}$ has to go up with $M_S$ to\navoid that the logarithm drives the correction unphysically large.\nThe optimal choice of $\\mu_\\mathrm{r}$, though, is somewhat empirical.}\nA good choice of the central scale $\\mu_0$ should come close to the\nstability point (plateau in the $\\mu_{r}$ variation) in the major part of THDM\nparameter space. Our choice \\eqref{eq:centralscale} of $\\mu_0$ \neffectively takes care of this and is eventually justified by the numerics.\n\nTo illustrate this and to estimate the theoretical uncertainties due to the\nresidual scale dependence, we compute the total width while the scale $\\mu_\\mathrm{r}$ is varied from $100{-}900\\unskip\\,\\mathrm{GeV}$.}\n{Results with central scale $\\mathswitch {M_\\Ph}$ are shown in App.~\\ref{app:mu0mh},\nproving that this would be not a good choice.}\nThe parameters $\\alpha$ and $\\beta$ are defined in the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme, and to compute results in other renormalization schemes their values are converted using\nEqs.~\\refeq{eq:convertinput1}--\\refeq{eq:convertinput3},\nwhich are solved numerically without linearization.\nThereafter the scale is varied, the RGEs solved, and the width computed using the respective renormalization scheme. \nThe results are shown \nin Fig.~\\ref{fig:plotmuscan} at LO (dashed) and NLO EW (solid) for the benchmark points Aa and Ab.\n\\begin{figure}\n \\centering\n \\subfigure[]{\n\\label{fig:plot_MUSCAN-Aa-L3MS}\n\\includegraphics{.\/diagrams\/plot_MUSCAN-Aa-L3MS.eps}\n}\n\\hspace{15pt}\n \\subfigure[]{\n\\label{fig:plot_MUSCAN-Ab-alphaMS}\n\\includegraphics{.\/diagrams\/plot_MUSCAN-Ab-L3MS.eps}\n}\n\\caption{The decay width for $\\mathswitchr h \\to 4f$ at LO \n(dashed) and NLO EW (solid) \nin dependence of the renormalization scale with $\\beta$ and $\\alpha$ defined in the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme. The result is computed in all four different renormalization schemes after converting the input at NLO (also for the LO curves)\nand displayed for the benchmark points Aa~(a) and Ab~(b) using the colour code of Fig.~\\ref{fig:plotrunningA}. }\n\\label{fig:plotmuscan}\n\\end{figure}\nThe QCD corrections are not part of the EW scale variation and therefore omitted in these results. The benchmark point Aa shows almost textbook-like behaviour with the LO computation \nexhibiting a strong scale dependence for all renormalization schemes, resulting in sizable differences between the curves. However, each of the NLO curves shows a wide extremum with a large plateau, reducing the scale dependence drastically, as it is expected for NLO calculations. \nThe central scale $\\mu_\\mathrm{r}=(\\mathswitch {M_\\Ph}+\\mathswitch {M_\\PH}+\\mathswitch {M_\\PAO}+2\\mathswitch {M_\\PHP})\/5$ lies perfectly in the middle of the plateau regions motivating this scale choice. In contrast, the naive scale choice $\\mu_0=\\mathswitch {M_\\Ph}$ is not within the plateau region, leads to large, unphysical corrections, and should not be chosen. The breakdown of the FJ($\\alpha)$ curve for small scales can be explained by the running which becomes unstable for these values (see Fig.~\\ref{fig:running_cba0.1-LM}).\nFor all renormalization schemes, the plateaus coincide and the agreement between the renormalization schemes is improved at NLO w.r.t.~the LO results. This is expected, since results obtained with different renormalization schemes should be equal up to higher-order terms, after the input parameters are properly converted. \nThe relative renormalization scheme dependence at the central scale,\n\\begin{align}\n\\Delta_\\mathrm{RS}=2 \\,\\frac{\\Gamma^{\\mathswitchr h \\to 4f}_\\mathrm{max}(\\mu_0)-\\Gamma^{\\mathswitchr h \\to 4f}_\\mathrm{min}(\\mu_0)}{\\Gamma^{\\mathswitchr h \\to 4f}_\\mathrm{max}(\\mu_0)+\\Gamma^{\\mathswitchr h \\to 4f}_\\mathrm{min}(\\mu_0)},\n\\label{eq:renschemedep}\n\\end{align}\n expresses the dependence of the result on the renormalization scheme. It can be computed from the difference of the smallest and largest width \nin the four renormalization schemes normalized to their average. \nIn the calculation of $\\Delta_\\mathrm{RS}$, the full NLO EW+QCD corrections to the width\n$\\Gamma^{\\mathswitchr h \\to 4f}$ should be taken into account.\nIn Tab.~\\ref{tab:schemevar}, $\\Delta_\\mathrm{RS}$ is given at LO and NLO and confirms the reduction of the scheme dependence in the NLO calculation. In addition, as already perceived when the running was analyzed, the $\\overline{\\mathrm{MS}}(\\alpha)$ scheme shows the smallest dependence on the renormalization scale, which attests a good absorption of further corrections into the NLO prediction.\\\\\n \\begin{table}\n \\centering\n \\renewcommand{\\arraystretch}{1.2}\n \\begin{tabular}{|c|cc|}\\hline\n \n & Scenario Aa & Scenario Ab\\\\\n $\\Delta^\\mathrm{LO}_\\mathrm{RS}$[\\%] & 0.67& 0.84 \\\\\n $\\Delta^\\mathrm{NLO}_\\mathrm{RS}$ [\\%] & 0.08 & 0.34 \\\\\\hline\n \\end{tabular} \n \\caption{The variation $\\Delta_\\mathrm{RS}$\nof the $\\mathswitchr h {\\to} 4f$ width using different renormalization schemes \nfor input parameters defined in the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme.}\n \\label{tab:schemevar}\n \\end{table}%\nThe situation for the benchmark point Ab is more subtle. For negative values of $c_{\\beta-\\alpha}$ the truncation of the schemes involving $\\lambda_3$ at \n$\\mu_\\mathrm{r}=250{-}300$ GeV as well as the breakdown of the running of the FJ($\\alpha)$ scheme, which both were observed in the running in Fig.~\\ref{fig:running_cba-0.1-LM}, are also ma\\-ni\\-fest in the computation of the $\\mathswitchr h {\\to} 4f$ width.\nTherefore, the results vary much more, and the extrema with the plateau regions are not as distinct as for the benchmark point Aa. They are even missing for the truncated curves. Nevertheless, the situation improves at NLO. \nAs for scenario~Aa, the central scale choice of $\\mu_0$ \nis more appropriate in contrast than the choice of $\\mathswitch {M_\\Ph}$.\n \nFor both benchmark points, \nthe estimate of the theoretical uncertainties by varying the scale by a factor of two from the central value for an arbitrary renormalization scheme is generally not appropriate. \nA proper strategy would be to identify the renormalization schemes which yield reliable results, and to use only those to quantify the theoretical uncertainties from the scale variation. In addition, the renormalization scheme dependence of those schemes should be investigated. \nThis procedure should be performed for different parameter regions (and corresponding \nbenchmark points) separately, which is beyond the scope of this work. \n \n\\subsection{\\boldmath{$c_{\\beta-\\alpha}$} dependence}\n\\label{sec:cbascanA}\n\nThe decay width for $\\mathswitchr h \\to 4f$ in dependence of $c_{\\beta-\\alpha}$ in scenario~A is presented in Fig.~\\ref{fig:plotscbascanA} for all renormalization schemes with \nthe input values $\\alpha$ and $\\beta$ defined in the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme. \n\\begin{figure}\n\\centering\n\\includegraphics{.\/diagrams\/plot_cbascan-A-diffschemes-L3MS.eps}\n\\caption{The decay width for $\\mathswitchr h \\to 4f$ at LO \n(dashed) and full NLO EW+QCD (solid) for scenario A in dependence of $c_{\\beta-\\alpha}$. \nThe input values are defined in the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme and \nare converted to the other schemes at NLO (also for the LO curves). \nThe results computed with different renormalization schemes are displayed with the colour code of Fig.~\\ref{fig:plotrunningA}, and the SM \n(with SM Higgs-boson mass $\\mathswitch {M_\\Ph}$) is shown for comparison in red.}\n\\label{fig:plotscbascanA}\n\\end{figure}%\nThe LO (dashed) and the full NLO EW+QCD total widths (solid) are computed in the different renormalization schemes after the NLO input conversion \n(without linearization) and using the constant default scale $\\mu_0$ of Eq.~\\eqref{eq:centralscale}. The SM values are illustrated in red. \nAt tree level the widths show the suppression w.r.t. to the SM with the factor $s_{\\beta-\\alpha}^2$ originating from the \n$\\mathswitchr H\\mathswitchr W\\PW$ and $\\mathswitchr H\\mathswitchr Z\\PZ$ couplings. The differences between the renormalization schemes are due to the conversion of the input. As the conversion induces NLO differences in the LO results, a pure LO computation is identical for all renormalization schemes as the conversion vanishes at this order \nand is represented by the LO curve of the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme. The suppression w.r.t.\\ the SM computation does not change at NLO, while the shape becomes slightly asymmetric,\nand the NLO results show a significantly better agreement between the renormalization schemes. \nDeviations of the THDM results from the SM expectations can be investigated when the SM Higgs-boson mass is identified with the mass \n$\\mathswitch {M_\\Ph}$ of the light CP-even Higgs boson~h \nof the THDM. The relative deviation of the full width from the SM is then\n\\begin{align}\n \\Delta_\\mathrm{SM}= \\frac{\\Gamma_\\mathrm{THDM}-\\Gamma_\\mathrm{SM}}{\\Gamma_\\mathrm{SM}},\n\\end{align}\nwhich is shown in Fig.~\\ref{fig:plot_cbascanrelSM-A-diffschemes-L3MS} at LO (dashed) and NLO (solid) in percent for parameters defined in the the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme.\n\\begin{figure}\n \\centering\n\\includegraphics{.\/diagrams\/plot_cbascanrelSM-A-diffschemes-L3MS.eps}\n \\caption{The relative difference of the decay width for $\\mathswitchr h \\to 4f$ in the THDM w.r.t.\\ the SM prediction at LO (dashed) and NLO EW+QCD (solid). \nThe input values are defined in the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme and \nare converted to the other schemes at NLO (also for the LO curves). \nThe results computed with different renormalization schemes are displayed with the colour code of Fig.~\\ref{fig:plotrunningA}.}\n\\label{fig:plot_cbascanrelSM-A-diffschemes-L3MS}\n\\end{figure}%\nThe SM exceeds the THDM widths at LO and NLO. \nThe LO shape which is just given by \n$c_{\\beta-\\alpha}^2$ shows minor distortions due to the parameter conversions. At NLO, the shape is slightly distorted by an asymmetry of the EW corrections, and a small offset of $-0.5$\\% is visible even in the alignment limit where the diagrams including heavy Higgs bosons still contribute. The NLO computations show larger negative deviations, and this could be used to improve current exclusion bounds or increase their significance. Nevertheless, in the whole scan region the deviation from the SM is within 6\\% and for phenomenologically most interesting region with $|c_{\\beta-\\alpha}|<0.1$ even less than 2\\%, which is challenging for experiments to measure.\n\n\\section{Conclusions}\n\\label{se:conclusion}\n\nConfronting experimental results on Higgs precision observables with theory predictions\nwithin extensions of the SM, provides an important alternative to search for\nphysics beyond the SM, in addition to the search for new particles.\nThe THDM comprises an extended scalar sector with regard to the SM Higgs sector and allows for a comprehensive study of the impact of new scalar degrees of freedom without introducing new fundamental symmetries or other new theoretical structures. \n\nIn this article, we have considered the Type I, II, lepton-specific, and flipped versions of the \nTHDM. We have introduced four different renormalization schemes which employ directly\nmeasurable parameters such as masses as far as possible and make use of fields that \ndirectly correspond to mass eigenstates.\nIn all the schemes, the masses are defined via on-shell conditions, the electric charge is fixed via the Thomson limit, and the coupling $\\lambda_5$ is defined with the $\\overline{\\mathrm{MS}}$ prescription. The fields are also defined on-shell which is most convenient in applications. The renormalization schemes differ in the treatment of the coupling $\\lambda_3$ and the mixing angles $\\alpha$ and $\\beta$: In the $\\overline{\\mathrm{MS}}(\\alpha)$ scheme, $\\alpha$ and $\\beta$ are renormalized using $\\overline{\\mathrm{MS}}$ conditions. In the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme instead $\\lambda_3$ and $\\beta$ are $\\overline{\\mathrm{MS}}$-renormalized parameters. In addition to the conventional treatment of tadpole contributions, we have\nimplemented an alternative prescription suggested by Fleischer and Jegerlehner where the mixing angles $\\alpha$ and $\\beta$ obtain extra terms of tadpole contributions, rendering these schemes gauge independent to all orders. \nIt should, however, be noted that the $\\overline{\\mathrm{MS}}(\\lambda_3)$ scheme is also gauge independent at NLO in the \nclass of $R_\\xi$ gauges. \nWe have also discussed relations to renormalization procedures\nsuggested in the literature for the THDM. \n\nA comparison of these four different renormalization schemes allows for testing \nthe perturbative consistency, and, for the parameter regions and renormalization schemes fulfilling this test, estimating the theoretical uncertainty due to the truncation of the perturbation series.\nTo further investigate the latter,\nwe have investigated the scale dependence, solving the corresponding RGEs. \nOne important observation is that it is crucial to be very careful and specific \nabout the definitions of the parameters applied, i.e.\\ the declaration of the renormalization\nscheme of the parameters is vital if one aims at precision.\nThis is already relevant in the formulation of\nbenchmark scenarios, because a conversion to a different scheme might alter the physical properties of the scenario significantly. For example, the alignment limit \nmay be reached with one specific set of parameters defined in a specific renormalization scheme, but converting these parameters consistently to parameters in a different renormalization scheme might shift the parameters away from the alignment limit.\n \nThe different renormalization schemes have been implemented into\na \\textsc{FeynArts} model file and are thus ready for applications.%\n\\footnote{The model file is restricted to a unit CKM matrix, but can be generalized\nto a non-trivial CKM matrix exactly as in the SM.}\nAs a first example, we have applied and tested the different schemes in the calculation of the decay width of a \nlight CP-even Higgs boson decaying into four massless fermions.\nWe discuss the dependence of the total $\\mathswitchr h{\\to}4f$ decay width on the renormalization scale\nand advocate a scale that is significantly higher than\nthe naive choice of $\\mu_\\mathrm{r} = \\mathswitch {M_\\Ph}$, taking care of the different mass scales in the\nTHDM Higgs sector.\nIn addition, results for various values of $\\cos (\\beta-\\alpha)$, a parameter entering the prediction already at LO, are presented. The deviations of the SM are relatively small, in the phenomenologically interesting region they are about \n$2{-}6\\%$---a challenge for future measurements. \n \nThe detailed\ndescription of the calculation of the decay width in the THDM \nand a survey of numerical results will be given in a forthcoming paper. \nThis includes a deeper investigation in the renormalization scale dependence and the comparison of different renormalization schemes \nfor more benchmark points\nas well as differential distributions.\n \n\n\\subsection*{Acknowledgements}\n\nWe would like to thank Ansgar Denner, Howard Haber and Jean-Nicolas Lang for helpful discussions\n{and especially Jean-Nicolas for an independent check of one-loop matrix elements\nagainst the crossing-related amplitudes used in \\citere{Denner:2016etu}.}\nHR's work is partially funded by the Danish National Research Foundation, grant number DNRF90. HR acknowledges also support by the National Science Foundation under Grant No. NSFPHY11-25915. \nWe thank the German Research Foundation~(DFG) and \nResearch Training Group GRK~2044 for the funding and the support and\nacknowledge support by the state of Baden--W\u00fcrttemberg through bwHPC and the \nDFG through grant no INST 39\/963-1 FUGG.\n\n\\section*{Appendix}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzcyma b/data_all_eng_slimpj/shuffled/split2/finalzzcyma new file mode 100644 index 0000000000000000000000000000000000000000..f54b5b826e499ef5116fe7e8d3a7e2977a177284 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzcyma @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\nAs early as in 1927~\\cite{dir27}, Paul A, M. Dirac considered the problem\nof extending Heisenberg's uncertainty relations to the\nLorentz-covariant world.\nIn 1945~\\cite{dir45}, he attempted to construct the Lorentz\ngroup using the Gaussian wave function.\nIn 1949~\\cite{dir49},\nDirac pointed out the task of constructing relativistic dynamics\nis to construct a representation of the inhomogeneous Lorentz group.\nHe then wrote down the ten generators of this group and their closed set\nof commutation relations. This set is known as the Lie algebra of\nthe Poincar\\'e group.\n\nIn 1963~\\cite{dir63}, Dirac considered two coupled harmonic oscillators\nand constructed an algebra leading to the Lie algebra for the $SO(3,2)$\nde Sitter group, which is the Lorentz group applicable to three space\ndimensions and two time-like variables.\n\nFrom the mathematical point of view, it is straight-forward to contract\none of those two time-like dimensions to construct $ISO(3,1)$ or the\nPoincar\\'e group. This is what we present in this paper. However,\nfrom the physical point of view, we are deriving the Poincar\\'e symmetry\nfor the Lorentz-covariant quantum world purely from the symmetries of\nHeisenberg's uncertainty relations.\n\nIn Sec.~\\ref{sp2}, it is noted that a one-dimensional uncertainty relation\ncontains the symmetry of the $Sp(2)$ group in the two-dimensional phase\nspace. It is pointed out that this group, with three generators, is isomorphic\nto the Lorentz group applicable to two space dimensions and one time variable.\nWe can next consider another set with three additional generators.\n\nIn Sec.~\\ref{2osc}, we write those Heisenberg uncertainty relations in\nterm of step-up and step-down operators in the oscillator system. It\nis then possible to consider the two coupled oscillator system\nwith the ten generators constructed by Dirac in 1963~\\cite{dir63}.\nIt is gratifying to note that this oscillator system can serve as the basic\nlanguage for the two-photon system of current interest~\\cite{yuen76,yurke86}.\n\nIn Sec.~\\ref{contraction}, we contract one of the time-like variables in\n$SO(3,2)$ to arrive at the inhomogeneous Lorentz group $ISO(3,1)$ or the\nPoincar\\'e group. In Sec.~\\ref{remarks}, we give some concluding\nremarks.\n\n\n\\section{Sp(2) Symmetry for the Single-variable Uncertainty Relation}\\label{sp2}\nIt is known that the symmetry of quantum mechanics and quantum field\ntheory is governed by the Poincar\\'e group~\\cite{dir49,dir62}. The\nPoincar\\'e group means the inhomogeneous Lorentz group which includes\nthe Lorentz group applicable to the four-dimensional Minkowskian\nspace-time, plus space-time translations~\\cite{wig39}.\n\nThe question is whether this Poincar\\'e symmetry is derivable from\nHeisenberg's uncertainty relation, which takes the familiar form\n\\begin{equation}\\label{101}\n \\left[x_{i}, p_{j}\\right] = i \\delta_{ij}.\n\\end{equation}\nThere are three commutation relations in this equation. Let us choose\none of them, and write it as\n\\begin{equation}\\label{102}\n [x, p] = i .\n\\end{equation}\nThis commutation relation possesses the symmetry of the Poisson bracket in\nclassical mechanics~\\cite{arnold78,stern84}. The best way to address this property\nis to use the Gaussian form for the Wigner function defined in the phase\nspace, which takes the\nform~\\cite{hkn88,kiwi90ajp,knp91}\n\\begin{equation}\\label{103}\nW(x,p) = \\frac{1}{\\pi} \\exp\\{-\\left(x^2 + p^2\\right)\\}.\n\\end{equation}\nThis distribution is concentrated in the circular region around the origin.\nLet us define the circle as\n\\begin{equation}\\label{105}\nx^{2} + p^{2} = 1.\n\\end {equation}\nWe can use the area of this circle in the phase space of $x$ and $p$ as the\nminimum uncertainty. This uncertainty is preserved\nunder rotations in the phase space:\n\\begin{equation}\\label{107}\n\\pmatrix{\\cos\\theta & -\\sin\\theta \\cr \\sin\\theta & \\cos\\theta }\n\\pmatrix{ x \\cr p },\n\\end{equation}\nas well as the squeeze of the form\n\\begin{equation}\\label{108}\n\\pmatrix{e^{\\eta} & 0 \\cr 0 & e^{-\\eta} }\n\\pmatrix{x \\cr p } .\n\\end{equation}\n\nThe rotation and the squeeze are generated by\n\\begin{equation}\\label{109}\nJ_{2} = - i\\left(x \\frac{\\partial}{\\partial p} - p\\frac{\\partial}{\\partial x } \\right),\n\\qquad\nK_{1} = -i \\left(x\\frac{\\partial}{\\partial x} - p \\frac{\\partial}{\\partial p} \\right),\n\\end{equation}\nrespectively.\nIf we take the commutation relation with these two operators, the result is\n\\begin{equation}\\label{111}\n\\left[J_{2}, K_{1}\\right] = -i K_{3},\n\\end{equation}\nwith\n\\begin{equation}\\label{113}\nK_{3} = -i \\left(x\\frac{\\partial}{\\partial p} +\np \\frac{\\partial}{\\partial x} \\right).\n\\end{equation}\nIndeed, these three generators form a closed set of commutation\nrelations:\n\\begin{equation}\\label{115}\n \\left[J_{2}, K_{1}\\right] = - iK_{3}, \\qquad\n \\left[J_{2}, K_{3}\\right] = iK_{1}, \\qquad\n \\left[K_{1}, K_{3}\\right] = iJ_{2}.\n\\end{equation}\nThis closed set is called the Lie algebra of the $Sp(2)$ group, isomorphic\nto the Lorentz group applicable to two space and one time dimensions.\n\nLet us consider the Minkowskian space of $(x, y, z, t)$. It is possible to\nwrite three four-by-four matrices satisfying the Lie algebra of Eq.(\\ref{115}).\nThe three four-by-four matrices satisfying this set\nof commutation relations are:\n\\begin{equation}\\label{119}\nJ_{2} = \\pmatrix{0 & 0 & i & 0 \\cr 0 & 0 & 0 & 0 \\cr i & 0 & 0 & 0 \\cr\n 0 & 0 & 0 & 0 },\\quad\nK_{1} = \\pmatrix{0 & 0 & 0 & i \\cr 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 \\cr\n i & 0 & 0 & 0 }, \\quad\n K_{3} = \\pmatrix{0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & i \\cr\n 0 & 0 & i & 0 } .\n\\end{equation}\nHowever, these matrices have null second rows and null second columns.\nThus, they can generate Lorentz transformations applicable only to the\nthree-dimensional space of $(x,z,t)$, while the $y$ variable remains invariant.\n\n\n\n\\section{Two-oscillator System}\\label{2osc}\nIn order to generate Lorentz transformations applicable to the full\nMinkowskian space, along with $J_{2},K_{1}$, and $K_{3}$ we need two\nmore Heisenberg commutation relations.\nIndeed, Paul A. M. Dirac started this program in 1963~\\cite{dir63}.\nIt is possible to write the two uncertainty relations using two harmonic\noscillators as\n\\begin{equation}\\label{121}\n\\left[a_{i}, a^{\\dag}_{j}\\right] = \\delta_{ij} .\n\\end{equation}\nwith\n\\begin{equation}\\label{123}\n a_{i} = \\frac{1}{\\sqrt{2}}\\left(x_{i} + ip_{i} \\right), \\qquad\n a^{\\dag}_{i} = \\frac{1}{\\sqrt{2}}\\left(x_{i} - ip_{i} \\right),\n\\end{equation}\nand\n\\begin{equation}\\label{125}\n x_{i} = \\frac{1}{\\sqrt{2}}\\left(a_{i} + a^{\\dag}_{i} \\right), \\qquad\n p_{i} = \\frac{i}{\\sqrt{2}}\\left(a^{\\dag}_{i} - a_{i} \\right),\n\\end{equation}\n where $i$ and $j$ could be 1 or 2.\n\nMore recently in 1986, this two-oscillator system was considered\nby Yurke {\\it et al.}~\\cite{yurke86} in their study of two-mode\ninterferometers. They considered first\n\\begin{equation}\\label{301}\n Q_{3} = \\frac{i}{2}\\left(a_{1}^{\\dag}a_{2}^{\\dag} - a_{1}a_{2}\\right),\n\\end{equation}\nwhich leads to the generation of the two-mode coherent state or the\nsqueezed state~\\cite{yuen76}.\n\nYurke {\\it et al.} then considered possible interferometers requiring\nthe following two additional operators.\n\\begin{equation}\\label{303}\nS_{3} = {1\\over 2}\\left(a^{\\dag}_{1}a_{1} + a_{2}a^{\\dag}_{2}\\right) ,\n\\qquad\nK_{3} = {1\\over 2}\\left(a^{\\dag}_{1}a^{\\dag}_{2} + a_{1}a_{2}\\right) .\n\\end{equation}\nThe three Hermitian operators from Eq.(\\ref{301}) and Eq.(\\ref{303})\nsatisfy the commutation relations\n\\begin{equation} \\label{305}\n\\left[K_{3}, Q_{3}\\right] = -iS_{3}, \\qquad\n\\left[Q_{3}, S_{3}\\right] = iK_{3}, \\qquad\n\\left[S_{3}, K_{3}\\right] = iQ_{3} .\n\\end{equation}\nThese relations are like those given in Eq.(\\ref{115}) for the Lorentz\ngroup applicable to two space-like and one time-like dimensions.\n\nIn addition, in the same paper~\\cite{yurke86}, Yurke {\\it et al.}\ndiscussed the possibility of constructing interferometers exhibiting\nthe symmetry generated by\n\\begin{equation}\\label{307}\n L_{1} = {1\\over 2}\\left(a^{\\dag}_{1}a_{2} + a^{\\dag}_{2}a_{1}\\right) , \\quad\n L_{2} = {1\\over 2i}\\left(a^{\\dag}_{1}a_{2} - a^{\\dag}_{2}a_{1}\\right), \\quad\n L_{3} = {1\\over 2}\\left(a^{\\dag}_{1}a_{1} - a^{\\dag}_{2}a_{2} \\right).\n\\end{equation}\nThese generators satisfy the closed set of commutation relations\n\\begin{equation}\\label{309}\n\\left[L_{i}, L_{j}\\right] = i\\epsilon_{ijk} L_{k} ,\n\\end{equation}\nand therefore define a Lie algebra which is the same as that for $SU(2)$\nor the three-dimensional rotation group.\n\nWe are then led to ask whether it is possible to construct a closed set\nof commutation relations with the six Hermitian operators from\nEqs.(\\ref{301},\\ref{303},\\ref{307}). It is not possible. We have to\nadd four additional operators, namely\n\\begin{eqnarray}\\label{311}\n&{}& K_{1} = -{1\\over 4}\\left(a^{\\dag}_{1}a^{\\dag}_{1} + a_{1}a_{1} -\n a^{\\dag}_{2}a^{\\dag}_{2} - a_{2}a_{2}\\right) , \\quad\nK_{2} = +{i\\over 4}\\left(a^{\\dag}_{1}a^{\\dag}_{1} - a_{1}a_{1} +\n a^{\\dag}_{2}a^{\\dag}_{2} - a_{2}a_{2}\\right) , \\nonumber \\\\[1ex]\n&{}& Q_{1} = -{i\\over 4}\\left(a^{\\dag}_{1}a^{\\dag}_{1} - a_{1}a_{1} -\n a^{\\dag}_{2}a^{\\dag}_{2} + a_{2}a_{2} \\right) , \\quad\nQ_{2} = -{1\\over 4}\\left(a^{\\dag}_{1}a^{\\dag}_{1} + a_{1}a_{1} +\n a^{\\dag}_{2}a^{\\dag}_{2} + a_{2}a_{2} \\right) .\n\\end{eqnarray}\n\nThere are now ten operators from Eqs.(\\ref{301},\\ref{303},\\ref{307},\\ref{311}).\nIndeed, these ten operators satisfy the following closed set of commutation\nrelations.\n\\begin{eqnarray}\\label{313}\n&{}& [L_{i}, L_{j}] = i\\epsilon _{ijk} L_{k} ,\\quad\n[L_{i}, K_{j}] = i\\epsilon_{ijk} K_{k} , \\quad\n[L_{i}, Q_{j}] = i\\epsilon_{ijk} Q_{k} , \\nonumber\\\\[1ex]\n &{}& [K_{i}, K_{j}] = [Q_{i}, Q_{j}] = -i\\epsilon _{ijk} L_{k} , \\quad\n [K_{i}, Q_{j}] = -i\\delta_{ij} S_{3}, \\nonumber\\\\[1ex]\n&{}&[L_{i}, S_{3}] = 0, \\quad [K_{i}, S_{3}] = -iQ_{i},\\quad\n[Q_{i}, S_{3}] = iK_{i} .\n\\end{eqnarray}\nAs Dirac noted in 1963~\\cite{dir63}, this set is the same as the Lie\nalgebra for the $SO(3,2)$ de Sitter group, with ten generators. This\nis the Lorentz group applicable to the three-dimensional space with\ntwo time variables. This group plays a very important role in\nspace-time symmetries.\n\n\nIn the same paper, Dirac pointed out that this set of commutation\nrelations serves as the Lie algebra for the four-dimensional\nsymplectic group commonly called $Sp(4)$, applicable to the systems\nof two one-dimensional particles, each with a two-dimensional phase\nspace.\n\n\nFor a dynamical system consisting of two pairs of canonical variables\n$x_{1}, p_{1}$ and $x_{2}, p_{2}$, we can use the four-dimensional\nspace with the coordinate variables defined as~\\cite{hkn95jmp}\n\\begin{equation}\\label{317}\n\\left(x_{1}, p_{1}, x_{2}, p_{2} \\right).\n\\end{equation}\nThen the four-by-four transformation matrix $M$ applicable to this\nfour-component vector is canonical if~\\cite{abra78,goldstein80}\n\\begin{equation}\\label{319}\nM J \\tilde{M} = J ,\n\\end{equation}\nwhere $\\tilde{M}$ is the transpose of the $M$ matrix,\nwith\n\\begin{equation}\\label{321}\nJ = \\pmatrix{0 & 1 & 0 & 0 \\cr -1 & 0 & 0 & 0 \\cr\n 0 & 0 & 0 & 1 \\cr 0 & 0 & -1 & 0 } .\n\\end{equation}\nAccording to this form of the $J$ matrix, the area of the phase space for\nthe $x_{1}$ and $p_{1}$ variables remains invariant, and the story is the\nsame for the phase space of $x_{2}$ and $p_{2}.$\n\nWe can then write the generators of the $Sp(4)$ group as\n\\begin{eqnarray}\\label{322}\n&{}& L_{1} = -\\frac{1}{2}\\pmatrix{0 & I\\cr I & 0 }\\sigma_{2} , \\quad\nL_{2} = \\frac{i}{2} \\pmatrix{0 & -I \\cr I & 0} I , \\quad\nL_{3} = \\frac{1}{2}\\pmatrix{-I & 0 \\cr 0 & I}\\sigma_{2} , \\nonumber\\\\[2ex]\n&{}& S_{3} = \\frac{1}{2}\\pmatrix{I & 0\\cr 0 & I} \\sigma_{2},\n\\end{eqnarray}\n\\noindent and\n\\begin{eqnarray}\\label{323}\n&{}&K_{1} = \\frac{i}{2}\\pmatrix{I & 0 \\cr 0 & -I } \\sigma_{1}, \\quad\nK_{2} = \\frac{i}{2} \\pmatrix{I & 0 \\cr 0 & I } \\sigma_{3}, \\quad\nK_{3} = -\\frac{i}{2}\\pmatrix{0 & I \\cr I & 0 } \\sigma_{1}, \\nonumber\\\\[2ex]\n&{}& Q_{1} = -\\frac{i}{2}\\pmatrix{I & 0 \\cr 0 & -I }\\sigma_{3}, \\quad\nQ_{2} = \\frac{i}{2}\\pmatrix{I & 0 \\cr 0 & I }\\sigma_{1} , \\quad\nQ_{3} = \\frac{i}{2}\\pmatrix{0 & I \\cr I & 0 } \\sigma_{3} ,\n\\end{eqnarray}\nwhere $I$ is the two-by-two identity matrix, while $\\sigma_{1},\n\\sigma_{2}$, and $\\sigma_{3}$ are the two-by-two Pauli matrices. The four matrices\ngiven in Eq.(\\ref{322}) generate rotations, while those of Eq.(\\ref{323}) lead to\nsqueezes in the four-dimensional phase space.\n\nAs for the difference in methods used in Secs.~\\ref{sp2} and~\\ref{2osc},\nlet us look at the ten four-by-four matrices given in Eqs.(\\ref{322}) and (\\ref{323}).\nAmong these ten matrices, six of them are diagonal. They are\n$S_{3}, L_{3}, K_{1}, K_{2} , Q_{1},$ and $Q_{2}$.\nIn the language of two harmonic oscillators, these generators do not\nmix up the first and second oscillators. There are six of them because\neach operator has three generators for its own $Sp(2)$ symmetry.\nLet us consider the three generators, $S_{3}, K_{2}$, and $Q_{2}$.\nFor each oscillator, the generators consist of\n\\begin{equation}\n \\sigma_{2}, \\quad i\\sigma_{1} \\quad \\mbox{and} \\quad i\\sigma_{3} .\n\\end{equation}\nThese separable generators thus constitute the Lie algebra of Sp(2) group\nfor the one-oscillator system, which we discussed in Sec.~\\ref{sp2}. Hence, the\none-oscillator system constitutes a subgroup of the two-oscillator system.\n\nThe off-diagonal matrix $L_{2}$ couples the first and second oscillators\nwithout changing the overall volume of the four-dimensional phase space.\nHowever, in order to construct the closed set of commutation relations,\nwe need the three additional generators: $L_{1} , K_{3},$ and $Q_{3}.$\nThe commutation relations given in Eq.(\\ref{313}) are clearly\nconsequences of Heisenberg's uncertainty relations.\n\n\n\n\n\n\n\\section{Contraction of SO(3,2) to ISO(3,1)}\\label{contraction}\n\n\nLet us next go back to the $SO(3,2)$ contents of this two-oscillator\nsystem~\\cite{dir63}. There are three space-like coordinates $(x, y, z)$\nand two time-like coordinates $s$ and $t$. It is thus possible to\nconstruct the five-dimensional space of $(x, y, z, t, s)$, and to consider\nfour-dimensional Minkowskian subspaces consisting of $(x, y, z, t)$\nand $(x, y, z, s)$.\n\nAs for the $s$ variable, we can make it longer or shorter, according\nto procedure of group contractions introduced first by In{\\\"o}n{\\\"u} and\nWigner~\\cite{inonu53}. In this five-dimensional space, the boosts along the\n$x$ direction with respect to the $t$ and $s$ variables are generated by\n\\begin{equation}\\label{333}\nA_{x} = \\pmatrix{0 & 0 & 0 & i & 0 \\cr 0 & 0 & 0 & 0 & 0 \\cr\n0 & 0 & 0 & 0 & 0 \\cr i & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0 }, \\qquad\nB_{x} = \\pmatrix{0 & 0 & 0 & 0 & i \\cr 0 & 0 & 0 & 0 & 0 \\cr\n0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0 \\cr i & 0 & 0 & 0 & 0 },\n\\end{equation}\nrespectively. The boost generators along the $y$ and $z$ directions\ntake similar forms.\n\nLet us then introduce the five-by-five contraction matrix~\\cite{kiwi87jmp,kiwi90jmp}\n\\begin{equation}\\label{335}\nC(\\epsilon) = \\pmatrix{1 & 0 & 0 & 0 & 0 \\cr 0 & 1 & 0 & 0 & 0 \\cr\n0 & 0 & 1 & 0 & 0 \\cr 0 & 0 & 0 & 1 & 0 \\cr 0 & 0 & 0 & 0 & \\epsilon } .\n\\end{equation}\nThis matrix leaves the first four columns and rows invariant, and the\nfour-dimensional Minkowskian sub-space of $(x, y, z, t)$ stays invariant.\n\nAs for the boost with respect to the $s$ variable, according to the procedure\nspelled out in Ref.~\\cite{kiwi87jmp,kiwi90jmp}, the contracted boost generator becomes\n\\begin{equation} \\label{341}\n B^{c}_{x} = \\lim_{\\epsilon\\rightarrow \\infty}\\frac{1}{\\epsilon}~\n \\left[C^{-1}(\\epsilon)~B_{x}~C(\\epsilon)\\right] =\n \\pmatrix{0 & 0 & 0 & 0 & i \\cr 0 & 0 & 0 & 0 & 0 \\cr\n0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0 } .\n\\end{equation}\nLikewise, $B^{c}_{y}$ and $B^{c}_{z}$ become\n\\begin{equation}\\label{343}\n B^{c}_{y} = \\pmatrix{0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & i \\cr\n 0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0 },\n \\qquad\n B^{c}_{z} = \\pmatrix{0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0 \\cr\n 0 & 0 & 0 & 0 & i \\cr 0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0} ,\n\\end{equation}\nrespectively.\n\nAs for the $t$ direction, the transformation applicable to the $s$ and $t$\nvariables is a rotation, generated by\n\\begin{equation}\\label{345}\nB_{t} = \\pmatrix{0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0 \\cr\n0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & i \\cr 0 & 0 & 0 & -i & 0 } .\n\\end{equation}\nThis matrix also becomes contracted to\n\\begin{equation}\\label{347}\nB^{c}_{t} = \\pmatrix{0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & 0 \\cr\n0 & 0 & 0 & 0 & 0 \\cr 0 & 0 & 0 & 0 & i \\cr 0 & 0 & 0 & 0 & 0 } .\n\\end{equation}\nThese contraction procedures are illustrated in Fig.~\\ref{contrac}.\n\n\n\n\\begin{figure\n\\centerline{\\includegraphics[scale=2.0]{contrac66.eps}}\n\\caption{Contraction of the $SO(3,2)$ group to the Poincar\\'e group.\nThe time-like $s$ coordinate is contracted with respect to the\nspace-like $x$ variable, and with respect to the time-like\nvariable $t$.}\\label{contrac}\n\\end{figure}\n\n\n\n\nThese four contracted generators lead to the five-by-five transformation\nmatrix\n\\begin{equation}\\label{349}\n\\exp\\left\\{-i\\left(aB^{c}_{x}+ bB^{c}_{y} + cB^{c}_{z} + dB^{c}_{t}\\right)\\right\\}\n= \\pmatrix{1 & 0 & 0 & 0 & a \\cr 0 & 1 & 0 & 0 & b \\cr\n0 & 0 & 1 & 0 & c \\cr 0 & 0 & 0 & 1 & d \\cr 0 & 0 & 0 & 0 & 1 } ,\n\\end{equation}\nperforming translations:\n\\begin{equation}\\label{351}\n \\pmatrix{1 & 0 & 0 & 0 & a \\cr 0 & 1 & 0 & 0 & b \\cr\n0 & 0 & 1 & 0 & c \\cr 0 & 0 & 0 & 1 & d \\cr 0 & 0 & 0 & 0 & 1 }\n\\pmatrix{x \\cr y \\cr z \\cr t \\cr 1 } =\n\\pmatrix{x + a \\cr y + b \\cr z + c \\cr t + d \\cr 1 } .\n\\end{equation}\nThis matrix leaves the first four rows and columns invariant. They are for\nthe Lorentz transformation applicable to the Minkowskian space of $(x, y, z, t)$.\n\nIn this way, the boosts along the $s$ direction become contracted to the\ntranslation. This means the group $SO(3,2)$ derivable from the Heisenberg's\nuncertainty relations becomes the inhomogeneous Lorentz group governing the\nPoincar\\'e symmetry for quantum mechanics and quantum field\ntheory in the Lorentz-covariant world~\\cite{dir63,dir62}.\n\nThe group contraction has a long history in physics, starting from the 1953\npaper by In{\\\"o}n{\\\"u} and Wigner~\\cite{inonu53}. It starts with a geometrical\nconcept. Our earth is a sphere, but is is convenient to consider a flat surface\ntangent to a given point on the spherical surface of the earth. This approximation\nis called the contraction of $SO(3)$ to $SE(2)$ or the two-dimensional Euclidean\ngroup with one rotational and two translational degrees of freedom.\n\nThis mathematical method was extended to the contraction of the $SO(3,1)$ Lorentz\ngroup to the three-dimensional Euclidean group. More recently, Kim and Wigner\nconsidered a cylindrical surface tangent to the sphere~\\cite{kiwi87jmp,kiwi90jmp}\nat its equatorial belt. This cylinder has one rotational degree of freedom and\none up-down translational degree of freedom. It was shown that the rotation and\ntranslation correspond to the helicity and gauge degrees of freedom for massless\nparticles.\n\nSince the Lorentz $SO(3,1)$ is isomorphic to the $SL(2,c)$ group of two-by-two\nmatrices, we can ask whether it is possible to perform the same contraction\nprocedure in the regime of two-by-two matrices. It does not appear possible\nto represent the $ISE(3)$ (inhomogeneous Euclidean group) with two-by-two\nmatrices. Likewise, there seem to be difficulties in addressing the\nquestion of contracting $SO(3,2)$ to $ISO(3,1)$ within the frame work of\nthe four-by-four matrices of $Sp(4)$.\n\n\n\n\n\n\\section{Concluding Remarks}\\label{remarks}\n\nSpecial relativity and quantum mechanics served as the major\ntheoretical basis for modern physics for one hundred years.\nThey coexisted in harmony: quantum mechanics augmented by Lorentz\ncovariance when needed.\nIndeed, there have been attempts in the past to construct a\nLorentz-covariant quantum world by augmenting the Lorentz group\nto the uncertainty relations~\\cite{dir27,dir45,dir49,dir62,yuka53}.\nThere are recent papers on this subject~\\cite{fkr71,lowe83,bars09}.\nThere are also papers on group contractions including contractions\nof the $SO(3,2)$ group~\\cite{gilmore74,bohm84,bohm85}.\n\nIt is about time for us to examine whether both of these two great\ntheories can be synthesized. The first step toward this process is\nto find the common mathematical ground. Before Newton, open orbits\n(comets) and closed orbits (planets) were treated differently, but\nNewton came up with one differential equation for both. Before\nMaxwell, electricity and magnetism were different branches of physics.\nMaxwell's equations synthesized these two branches into one. It is\nshown in this paper that the group $ISO(3,1)$ can be derived from\nthe algebra of quantum mechanics.\n\nIt is gratifying to note that the Poincar\\'e symmetry is derivable\nwithin the system of Heisenberg's uncertainty relations. The\nprocedure included two coupled oscillators resulting in the $SO(3,2)$\nsymmetry~\\cite{dir63}, and the contraction of this $SO(3,2)$ to the\ninhomogeneous Lorentz group $ISO(3,1)$.\n\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\n\n\n\\subsection{Experimental Procedure}\nWe study four experimental factors: the model, the batch size, the initialization point, and the learning rate. The model factor has three levels that correspond to the three models described above. The batch size will take the values corresponding to SGD-$1$, SGD-$200$, SGD-$500$, and Gradient Descent. The initialization point will be randomly selected from a uniform distribution on a ball whose radius will take values $\\lbrace 10^{-3},10^{-2},10^{-1}, 1 \\rbrace$ and is centered at either the sharpest minimizer or the flattest minimizer. The learning rates will be linear combinations of the upper and lower bounds. Specifically, if $u$ is the threshold for convergence and $l$ is the threshold for divergence, then the learning rates will be $1.5l$, $0.5(u+l)$, or $0.5u$. For each unique collection of the levels of the four factors, one hundred independent runs are executed with at most twenty iterations. For each run, the euclidean distances between the iterates and the minimizer near the initialization point are recorded. \n\n\\subsection{Results and Discussion}\nNote that the results Models 1, 2 and 3 are nearly identical for the purposes of our discussion, and so we will feature Model 1 only in our discussion below.\n\n\\begin{figure}[hbtp]\n\\centering \n\\includegraphics[width=\\textwidth]{st_fig1}\n\\caption[SGD-$k$ on Styblinski-Tang Near Flat Minimizer]{The behavior of SGD-$k$ on Model 1 for different batch sizes when initialized near the flat minimizer. The $y$-axis shows the distance (in logarithmic scale) between the iterates and the flat minimizer for all runs of the specified batch size and specified learning rate.}\n\\label{figure-sgd:st-1}\n\\end{figure}\n\n\\begin{figure}[htbp]\n\\centering \n\\includegraphics[width=\\textwidth]{st_fig3}\n\\caption[SGD-$k$ on Styblinski-Tang Near Sharp Minimizer]{The behavior of SGD-$k$ on Model 1 for different batch sizes when initialized near the sharp minimizer. The $y$-axis shows the distance (in logarithmic scale) between the iterates and the sharp minimizer for all runs of the specified batch size and specified learning rate.}\n\\label{figure-sgd:st-3}\n\\end{figure}\n\n\\begin{figure}[htbp]\n\\centering\n\\includegraphics[width=\\textwidth]{st_fig2}\n\\caption[Initialization and SGD-$k$ on Styblinski-Tang Near Flat Minimizer]{The behavior of SGD-$k$ on Model 1 for different starting radii when initialized near the flat minimizer. The $y$-axis shows the distance (in logarithmic scale) between the iterates and the flat minimizer for all runs of the specified starting radius and specified learning rate.}\n\\label{figure-sgd:st-2}\n\\end{figure}\n\n\\begin{figure}[hbtp]\n\\centering \n\\includegraphics[width=\\textwidth]{st_fig4}\n\\caption[Initialization and SGD-$k$ on Styblinski-Tang Near Sharp Minimizer]{The behavior of SGD-$k$ on Model 1 for different starting radii when initialized near the sharp minimizer. The $y$-axis shows the distance (in logarithmic scale) between the iterates and the sharp minimizer for all runs of the specified starting radius and specified learning rate.}\n\\label{figure-sgd:st-4}\n\\end{figure}\n\nFigure \\ref{figure-sgd:st-1} shows the distance between SGD-$k$ iterates and the flat minimizer on Model 1 for different batch sizes and different learning rates when SGD-$k$ is initialized near the flat minimizer. \nWe note that, regardless of batch size, the iterates are diverging from the flat minimizer and are converging to the corners of the feasible region for learning rates $1.5l$ and $0.5(u+l)$. \nOn the other hand, for the learning rate $0.5u$, we see stability for $k=1,200,500$ about the minimizer and we see convergence for GD (i.e., $k=\\infty$). \nSimilarly, Figure \\ref{figure-sgd:st-3} shows the distance between SGD-$k$ iterates and the sharp minimizer on Model 1 for different batch sizes and different learning rates when SGD-$k$ is initialized near the sharp minimizer. \nWe note that, regardless of batch size, the iterates are diverging from the sharp minimizer and converging to the corners of the feasible region for learning rates $1.5l$ and $0.5(u+l)$. On the other hand, for the learning rate $0.5u$, we see stability for $k=1,200,500$ about the sharp minimizer and we see convergence for GD (i.e., $k = \\infty$). \n\nTaking the results in Figures \\ref{figure-sgd:st-1} and \\ref{figure-sgd:st-3} together, we see that we are able to use our deterministic mechanism to find learning rates that either ensure divergence from or stability about a minimizer if we know its local geometric properties. Consequently, we again have evidence that our deterministic mechanism can correctly predict the behavior of SGD-$k$ for nonconvex problems. Moreover, when divergence does occur, Figures \\ref{figure-sgd:st-1} and \\ref{figure-sgd:st-3} also display exponential divergence as predicted by our deterministic mechanism.\n\nFor a different perspective, Figure \\ref{figure-sgd:st-2} shows the distance between SGD-$k$ iterates and the flat minimizer on Model 1 for different starting radii and different learning rates when SGD-$k$ is initialized near the flat minimizer. We note that, regardless of the starting radius, the iterate are diverging from the flat minimizer and are converging to the corners of the feasible region for learning rates $1.5l$ and $0.5(u+l)$. On the other hand, for the learning rate $0.5u$, we see stability and even convergence for all of the runs regardless of the starting radius. \nSimilarly, Figure \\ref{figure-sgd:st-4} shows the distance between SGD-$k$ iterates and the sharp minimizer on Model 1 for different starting radii and different learning rates when SGD-$k$ is initialized near the sharp minimizer. We note that, regardless of the starting radius, the iterates are diverging from the sharp minimizer and are converging to the corners of the feasible region for learning rates $1.5l$ and $0.5(u+l)$. On the other hand, for learning rate $0.5u$, we see stability about and even convergence to the sharp minimizer. \n\nTaking the results in Figures \\ref{figure-sgd:st-2} and \\ref{figure-sgd:st-4} together, we see that we are able to use our deterministic mechanism to find learning rates that either ensure divergence from, or stability about, a minimizer if we know its local geometric properties. Consequently, we again have evidence that our deterministic mechanism can correctly predict the behavior of SGD-$k$ for nonconvex problems. Moreover, when divergence does occur, Figures \\ref{figure-sgd:st-1} and \\ref{figure-sgd:st-3} also display exponential divergence as predicted by our deterministic mechanism.\n\n\\subsection{Experimental Procedure} \nWe study four experimental factors: the model, the learning method, the initialization point, and the learning rate. \nThe model factor has four levels as given by the four model objective functions shown in Figure \\ref{figure-sgd:qc-models}. The learning method has two levels which are SGD-$1$ or GD. \nThe initialization point also has two levels: either the initialization point is selected randomly with a uniform probability from a disc of radius $10^{-8}$ centered about the circular basin minimizer, or the initialization point is selected randomly with a uniform probability from a disc of radius $10^{-8}$ centered about the quadratic basin minimizer. The learning rate has levels that are conditional on the initialization point. \nIf the initialization point is near the minimizer of the circular basin, then the learning rate takes on values $10^{10}$, $5(10^{10})$, $10^{11}$, $5(10^{11})$, $10^{12}$, and $5(10^{12})$. \nIf the initialization point is near the minimizer of the circular basin, then the learning rate takes on values $1, 4, 16, 64, 256, 1024$. \n\nFor each unique collection of the levels of the four factors, one hundred independent runs are executed with at most twenty iterations. For each run, the euclidean distance between the iterates and the circular basin minimizer and the euclidean distance between the iterates and the quadratic basin minimizer are recorded. \n\n\\subsection{Results and Discussion}\nNote that the results of SGD-$k$ on Models 1 and 2 and Models 3 and 4 are similar when initialized around the circular basin minimizer, and the results of SGD-$k$ on Models 1 and 3 and Models 2 and 4 are similar when initialized around the quadratic basin minimizer. Hence, when we discuss the circular basin minimizer, we will compare Models 1 and 3, but we could have just as easily replaced Model 1 with Model 2 or Model 3 with Model 4 and the discussion would be identical. Similarly, when we discuss the quadratic basin minimizer, we will compare Models 1 and 2, but we could have replaced the results of Model 1 with Model 3 or Model 2 with Model 4 and the discussion would be identical.\n\n\\begin{figure}[hbt]\n\\centering\n\\includegraphics[width=\\textwidth]{cq_fig1}\n\\caption[SGD-$1$ Near Circle Basin Minimizer]{The behavior of SGD-$1$ on Models 1 and 3 when initialized near the circular minimum. The $y$-axis shows the distance (in logarithmic scale) between the estimates and the circular minimum for all runs of the specified model and the specified learning rate.}\n\\label{figure-sgd:cq-1}\n\\end{figure}\n\n\\begin{figure}[hbt]\n\\centering\n\\includegraphics[width=\\textwidth]{cq_fig2}\n\\caption[SGD-$1$ Near Quadratic Basin Minimizer]{The behavior of SGD-$1$ on Models 1 and 2 when initialized near the quadratic minimum. The $y$-axis shows the distance (in logarithmic scale) between the iterates and the quadratic minimum for all runs of the specified model and the specified learning rate.}\n\\label{figure-sgd:cq-2}\n\\end{figure}\n\nFigure \\ref{figure-sgd:cq-1} shows the distance between the iterates and the circular basin minimizer for the one hundred independent runs of SGD-$1$ for the specified model and the specified learning rate when initialized near the circular basin minimizer. The learning rates that are displayed are the ones where a transition in the convergence-divergence behavior of the method occur for the specific model. Specifically, SGD-$1$ begins to diverge for learning rates between $5 \\times 10^{10}$ and $10^{11}$ for Model 1 and between $5\\times 10^{11}$ and $10^{12}$ for Model 3. \nSimilarly, Figure \\ref{figure-sgd:cq-2} shows the distance between the iterates and the quadratic basin minimizer for the one hundred independent runs of SGD-$1$ for the specified model and the specified learning rate when initialized near the quadratic basin minimizer. SGD-$1$ begins to diverge for learning rates between $4$ and $16$ for Model 1 and between $256$ and $1024$ for Model 2. \n\nFrom these observations, we see that relatively flatter minimizers enjoy larger thresholds for divergence of SGD-$1$ in comparison to sharper minimizers. Moreover, while the bounds computed in Table \\ref{table:lb} are conservative (as we would expect), they are still rather informative, especially in the case of the quadratic basin. Thus, regarding questions (1) and (2) above, we see that while our thresholds are slightly conservative, we still correctly predict the expected behavior of the SGD-$k$ iterates. Moreover, regarding question (3), we observe that the iterates diverge from their respective minimizers in a deterministic way with an exponential rate. Again, this observations lends credence to our deterministic mechanism over the widely accepted stochastic mechanism.\n\n\n\\begin{figure}[hbt]\n\\centering\n\\includegraphics[width=\\textwidth]{cq_fig3}\n\\caption[GD Near Circle Basin Minimizer]{The behavior of GD on Models 1 and 3 when initialized near the circular basin minimizer. The $y$-axis shows the distance (in logarithmic scale) between the iterates and the circular minimizer for all runs of the specified model and the specified learning rate.}\n\\label{figure-sgd:cq-3}\n\\end{figure}\n\nFigure \\ref{figure-sgd:cq-3} shows the distance between the iterates and the circular basin minimizer for the one hundred independent runs of gradient descent (GD) for the specified model and the specified learning rate when initialized near the circular basin minimizer. If we compare Figures \\ref{figure-sgd:cq-1} and \\ref{figure-sgd:cq-3}, we notice that, for Model 1 and learning rate $5 \\times 10^{10}$, the runs of GD converge whereas some of the runs for SGD-$1$ diverge. Although we do not report the results for all of the models or all of the learning rates, we note the boundary for divergence-convergence for GD are smaller than those of SGD-$1$. In light of our deterministic mechanism, this behavior is expected: as the batch-size, $k$, increases, the lower bound on the learning rates for divergence increases. Therefore, we should expect that at those boundary learning rates where some SGD-$1$ runs are stable and others diverge, for GD, we should only see stable runs, and, indeed, this is what the comparison of Figures \\ref{figure-sgd:cq-1} and \\ref{figure-sgd:cq-3} shows.\n\n\\begin{figure}[hbt]\n\\centering\n\\includegraphics[width=\\textwidth]{cq_fig4}\n\\caption[Comparing SGD-$1$ and GD Near Quadratic Basin Minimizer]{The behavior of SGD-$1$ and GD on Model 1 when initialized near the quadratic minimum for select learning rates. The $y$-axis shows the distance (in logarithmic scale) between the iterates and the circular minimizer for all runs of the specified method and the specified learning rate.}\n\\label{figure-sgd:cq-4}\n\\end{figure}\n\nFigure \\ref{figure-sgd:cq-4} shows the distance between the iterates and the circular basin minimizer for the one hundred independent runs of SGD-$1$ and the one hundred independent runs of GD for the specified method and the specified learning rate when initialized near the quadratic basin minimizer. In Figure \\ref{figure-sgd:cq-4}, we see that for the learning rates that lead to divergence from the quadratic basin minimizer (compare to the top three subplots of Figure \\ref{figure-sgd:cq-1}) are in some cases able to converge to the circular minimizer. Again, in light of our deterministic mechanism, this is expected: the learning rates shown in Figure \\ref{figure-sgd:cq-4} guarantee divergence from the quadratic minimizer and are sufficiently small that they can lead to convergence to the circular minimizer. However, we notice in Figure \\ref{figure-sgd:cq-4} that as the learning rate increases, even though the learning rate is below the divergence bound for the circular minimizer, the iterates for SGD-$1$ and GD are diverging from both the circular and quadratic minima and are converging to the corners of the feasible region.\n\nTo summarize, we see that convergence-divergence behavior of SGD-$k$ for the quadratic-circle sums nonconvex problem is captured by our deterministic mechanism. Although our estimated bounds are sometimes conservative, they generally divide the convergence and divergence regions of SGD-$k$. Moreover, as our deterministic mechanism predicts, we observed exponential divergence away from minimizers when the learning rate is above the divergence threshold. Again, this lends credence to our mechanism for divergence and provides evidence against the stochastic mechanism for SGD-$k$ iterates to ``escape'' sharp minimizers.\n\n\\subsection{A Generic Stochastic Quadratic Problem} \\label{subsection:stochastic-quadratic}\n\nWe will begin by defining the quadratic problem, derive some of its relevant properties, and discuss the two types of minimizers that can occur for the problem.\n\n\\begin{problem}[Quadratic Problem] \\label{problem-sgd:quadratic}\nLet $(\\Omega,\\mathcal{F},\\mathbb{P})$ be a probability space. Let $Q \\in \\mathbb{R}^{p \\times p}$ be a nonzero, symmetric, positive semidefinite random matrix and $r \\in \\mathbb{R}^p$ be a random vector in the image space of $Q$ (with probability one) such that (1) $\\E{Q} \\prec \\infty$, (2) $\\E{ Q \\E{Q} Q} \\prec \\infty$, (3) $\\E{Q \\E{Q} r} $ is finite, and (4) $\\E{r'Qr}$ is finite. Let $\\lbrace (Q_N,r_N): N \\in \\mathbb{N} \\rbrace$ be independent copies of $(Q,r)$. The quadratic problem is to use $\\lbrace (Q_N,r_N) : N \\in \\mathbb{N} \\rbrace$ to determine a $\\theta^*$ such that\n\\begin{equation} \\label{eqn-sgd-problem:quadratic-obj}\n\\theta^* \\in \\underset{\\theta \\in \\mathbb{R}^p}{\\mathrm{argmin}} \\frac{1}{2}\\theta'\\E{Q}\\theta + \\E{r}'\\theta.\n\\end{equation}\n\\end{problem}\n\nThere are several aspects of the formulation of Problem \\ref{problem-sgd:quadratic} worth highlighting. First, $Q$ is required to be symmetric and positive semidefinite, and $r$ is required to be in the image space of $Q$. Together, these requirements on $(Q,r)$ ensure that for each $\\omega$ in some probability one set in $\\mathcal{F}$, there exists a minimizer for $0.5 \\theta'Q\\theta + r'\\theta$. While this condition seems quite natural for the quadratic problem, it will cause some challenges for generalizing our results directly to the nonconvex case, which we will discuss further in \\S \\ref{subsection:deterministic-mechanism}. Second, we have required a series of assumptions on the moments of $(Q,r)$. As we will see in the analysis of SGD-$k$ on Problem \\ref{problem-sgd:quadratic}, these moment assumptions will be essential to the stability of the iterates. Finally, we have not required that the solution, $\\theta^*$, be unique. That is, we are not requiring a strongly convex problem. This generality is important as it will allow us to better approximate more exotic nonconvex problems. \n\nWe now establish some basic geometric properties about the quadratic problem. We first establish a standard, more convenient reformulation of the objective function in (\\ref{eqn-sgd-problem:quadratic-obj}). Then, we formulate some important geometric properties of the quadratic problem.\n\n\\begin{lemma} \\label{lemma-sgd:objective-reformulation}\nThe objective function in (\\ref{eqn-sgd-problem:quadratic-obj}) is equivalent to (up to an additive constant)\n\\begin{equation}\n\\frac{1}{2}(\\theta - \\theta^*)'\\E{Q}(\\theta - \\theta^*),\n\\end{equation}\nfor any $\\theta^*$ satisfying (\\ref{eqn-sgd-problem:quadratic-obj}). \n\\end{lemma}\n\\begin{proof}\nIf $\\theta^*$ is a minimizer of the objective in (\\ref{eqn-sgd-problem:quadratic-obj}), then $\\theta^*$ must be a stationary point of the objective function. This implies that $-\\E{Q}\\theta^* = \\E{r}$. Therefore, the objective in (\\ref{eqn-sgd-problem:quadratic-obj}) is, up to an additive constant,\n\\begin{equation}\n\\frac{1}{2} \\theta'\\E{Q}\\theta - \\theta' \\E{Q} \\theta^* + \\frac{1}{2}(\\theta^*)'\\E{Q} \\theta^* = \\frac{1}{2}(\\theta - \\theta^*)'\\E{Q}(\\theta - \\theta^*),\n\\end{equation} \n\\end{proof}\n\nBecause $Q$ is symmetric, positive semidefinite with probability one, $\\E{Q}$ will also be symmetric, positive semidefinite. That is, $m := \\texttt{rank}(\\E{Q}) \\leq p$. For future reference, we will denote the nonzero eigenvalues of $\\E{Q}$ by $\\lambda_1 \\geq \\cdots \\lambda_m > 0$. Beyond the eigenvalues, we will also need some higher order curvature information, namely, $s_Q$ and $t_Q$, which are given by\n\\begin{equation} \\label{eqn-sgd:tQ}\nt_Q = \\sup \\left\\lbrace \\frac{v'\\E{Q \\E{Q} Q}v - v' \\E{Q}^3 v}{v' \\E{Q} v} : v \\in \\mathbb{R}^p, ~ v'\\E{Q}v \\neq 0 \\right\\rbrace,\n\\end{equation}\nand\n\\begin{equation} \\label{eqn-sgd:sQ}\ns_Q = \\inf \\left\\lbrace \\frac{v'\\E{Q \\E{Q} Q}v - v' \\E{Q}^3 v}{v' \\E{Q} v} : v \\in \\mathbb{R}^p, ~ v'\\E{Q}v \\neq 0 \\right\\rbrace.\n\\end{equation}\nFrom the definitions in (\\ref{eqn-sgd:tQ}) and (\\ref{eqn-sgd:sQ}), it follows that $t_Q \\geq s_Q$. The next result establishes that for nontrivial problems, $s_Q > 0$. \n\n\\begin{lemma} \\label{lemma-sgd:higher-moment-bounds}\nLet $Q$ be as in (\\ref{problem-sgd:quadratic}). Then, the following properties hold.\n\\begin{enumerate}\n\\item If there is an $x \\in \\mathbb{R}^p$ such that $x' \\E{Q} x = 0$ then $x' \\E{ Q \\E{Q} Q} x = 0$. \n\\item $\\E{Q \\E{Q} Q} \\succeq \\E{Q}^3$. \n\\item If $\\Prb{Q = \\E{Q}} < 1$ then $s_Q > 0$. \n\\end{enumerate} \n\\end{lemma}\n\\begin{proof}\nFor (1), when $x=0$ the result follows trivially. Suppose then that there is an $x \\neq 0$ such that $x'\\E{Q}x = 0$. Then, since $Q \\succeq 0$, $x'Qx = 0$ almost surely. Therefore, $x$ is in the null space of all $Q$ on a probability one set. Therefore, $x'Q \\E{Q} Q x = 0$ with probability one. The conclusion follows.\n\nFor (2), note that for any $x \\in \\mathbb{R}^q$, the function $f(M) = x'M \\E{Q} Mx$ is convex over the space of all symmetric, positive semi-definite matrices. Therefore, by Jensen's inequality, $\\E{f(Q)} \\geq f(\\E{Q})$. Moreover, since $x$ is arbitrary, the conclusion follows.\n\nFor (3), note that Jensen's inequality holds with equality only if $f$ is affine or if $\\Prb{Q = \\E{Q}} = 1$. Therefore, if $\\Prb{Q = \\E{Q}} < 1$, then the second condition is ruled out. Moreover, if $\\Prb{Q = \\E{Q}} < 1$, then $\\Prb{Q = 0} < 1$. Thus, $\\exists x \\neq 0$ such that $f(Q) \\neq 0$, which rules out the first condition. \n\\end{proof}\n\nNow, if $\\Prb{Q = \\E{Q}} = 1$ then there exists a solution to (\\ref{eqn-sgd-problem:quadratic-obj}) such that $Q\\theta^* + r = 0$ with probability one. This phenomenon of the existence of a single solution for all pairs $(Q,r)$ with probability one will play a special role, and motivates the following definition.\n\n\\begin{definition}[Homogeneous and Inhomogeneous Solutions]\nA solution to $\\theta^*$ satisfying (\\ref{eqn-sgd-problem:quadratic-obj}) is called homogeneous if\n\\begin{equation} \\label{eqn-sgd-definition:homogeneous}\n\\theta^* \\in \\mathrm{argmin}_{\\theta \\in \\mathbb{R}^p} \\frac{1}{2} \\theta'Q \\theta + r'\\theta \\text{ with probability one.}\n\\end{equation}\nOtherwise, the solution is called inhomogeneous.\n\\end{definition}\n\nIn the case of a homogeneous minimizer, the objective function in (\\ref{eqn-sgd-definition:homogeneous}) can be rewritten as (up to an additive constant)\n\\begin{equation}\n\\frac{1}{2} (\\theta - \\theta^*)' Q (\\theta - \\theta^*),\n\\end{equation}\nwhich follows from the same reasoning Lemma \\ref{lemma-sgd:objective-reformulation}. Importantly, SGD-$k$ will behave differently for problems with homogeneous minimizers and inhomogeneous minimizer. As we will show, for homogeneous minimizers, SGD-$k$ behaves rather similarly to classical gradient descent in terms of convergence and divergence, whereas for inhomogeneous minimizers, SGD-$k$ will behave markedly differently. To show these results, we will first define SGD-$k$ for the quadratic problem.\n\n\\begin{definition}[SGD-$k$] \\label{definition-sgd:sgd-k}\nLet $\\theta_0 \\in \\mathbb{R}^p$ be arbitrary. For Problem \\ref{problem-sgd:quadratic}, SGD-$k$ generates a sequence of iterates $\\lbrace \\theta_N : N \\in \\mathbb{N} \\rbrace$ defined by\n\\begin{equation} \\label{eqn-definition:sgd-k}\n\\theta_{N+1} = \\theta_N - \\frac{C_{N+1}}{k}\\sum_{j=Nk+1}^{N(k+1)} Q_j\\theta_N + r_j,\n\\end{equation} \nwhere $\\lbrace C_N: N \\in \\mathbb{N} \\rbrace$ is a sequence of scalars. \n\\end{definition}\n\nIn Definition \\ref{definition-sgd:sgd-k}, when $k = 1$, we have the usual notion of stochastic gradient descent and, as $k \\to \\infty$, we recover gradient descent. Therefore, this notion of SGD-$k$ will allow us to generate results that explore the complete range of stochastic to deterministic methods, and bridge the theory between stochastic gradient descent and gradient descent.\n\n\\subsection{SGD-$k$ on the Stochastic Quadratic Problem} \\label{subsection:analysis}\n\n\\begin{table}\n\\tbl{Summary of Notation}\n{\n\\begin{tabular}{@{}ll@{}} \\toprule\nNotation & Value \\\\ \\midrule\n$M$ & $\\E{Q \\E{Q} Q} - \\E{Q}^3$ \\\\\n$e_N$ & $(\\theta_N - \\theta^*)'\\E{Q}(\\theta_N - \\theta^*)$ \\\\\n$e_{N,l}$ & $(\\theta_N - \\theta^*)'\\E{Q}^l(\\theta_N - \\theta^*)$ for $l \\in \\mathbb{N}_{\\geq 2}$ \\\\\n$e_{N,M}$ & $(\\theta_N - \\theta^*)'M (\\theta_N - \\theta^*)$ \\\\ \n$m$ & $\\texttt{rank}( \\E{Q} )$ \\\\\n$\\lambda_1 \\geq \\cdots \\geq \\lambda_m > 0$ & Nonzero eigenvalues of $\\E{Q}$ \\\\ \n$t_Q, s_Q$ & Higher-order Curvature Parameters of $Q$ \\\\\\bottomrule\n\\end{tabular}\n}\n\\label{table:summary-of-notation}\n\\end{table}\n\nTable \\ref{table:summary-of-notation} summarizes the notation that will be used throughout the remainder of this work. With this notation, we now analyze the iterates generated by SGD-$k$ for Problem \\ref{problem-sgd:quadratic}. Our first step is to establish a relationship between $e_{N+1}$ and $e_N$, which is a standard computation using conditional expectations.\n\n\\begin{lemma} \\label{lemma-sgd:quadratic-recursion}\nLet $\\lbrace \\theta_N : N \\in \\mathbb{N} \\rbrace$ be the iterates of SGD-$k$ from Definition \\ref{definition-sgd:sgd-k}. Let $\\theta^*$ be a solution to the quadratic problem. Then, with probability one, \n\\begin{equation}\n\\begin{aligned}\n\\cond{e_{N+1}}{\\theta_N} &= e_{N} - 2C_{N+1}e_{N,2} + C_{N+1}^2 e_{N,3} + \\frac{C_{N+1}^2}{k} e_{N,M} \\\\\n\t\t\t\t\t\t&+ 2\\frac{C_{N+1}^2}{k}(\\theta_N - \\theta^*)' \\E{ Q \\E{Q} (Q\\theta^* + r)} \\\\\n\t\t\t\t\t\t&+ \\frac{C_{N+1}^2}{k} \\E{ (Q\\theta^* + r)' \\E{Q}(Q \\theta^* + r)}.\n\\end{aligned}\n\\end{equation}\nMoreover, if $\\theta^*$ is a homogeneous minimizer, then, with probability one,\n\\begin{equation}\n\\begin{aligned}\n\\cond{e_{N+1}}{\\theta_N} &= e_{N} - 2C_{N+1} e_{N,2} + C_{N+1}^2 e_{N,3} + \\frac{C_{N+1}^2}{k} e_{N,M}.\n\\end{aligned}\n\\end{equation}\n\\end{lemma}\n\\begin{proof}\nThe result is a straightforward calculation from the properties of SGD-$k$ and the quadratic problem. In the case of the homogeneous minimizer, recall that $Q\\theta^* +r = 0$ with probability one.\n\\end{proof} \n\nOur second step is to find bounds on $\\cond{e_{N+1}}{\\theta_N}$ in terms of $e_N$. While this is a standard thing to do, our approach differs from the standard approach in two ways. First, not only will we find an upper bound, but we will also find a lower bound. Second, our lower bounds will be much more delicate than the usual strong convexity and Lipschitz gradient based arguments \\cite[see][\\S 4]{bottou2016}. This second step can be accomplished for the homogeneous case with the following two lemmas.\n\n\\begin{lemma} \\label{lemma-sgd:technical-systematic-component}\nIn the setting of Lemma \\ref{lemma-sgd:quadratic-recursion} and for $\\alpha_j \\geq 0$ for $j=0,1,2,3$, \n\\begin{equation} \\label{eqn-sgd-lemma:technical-systematic-component}\n\\begin{aligned}\n&\\alpha_0e_N - \\alpha_1 C_{N+1} e_{N,2} + \\alpha_2 C_{N+1}^2 e_{N,3} + \\alpha_3 C_{N+1}^2 e_{N,M}\n\\end{aligned}\n\\end{equation}\nis bounded above by\n\\begin{equation}\ne_N\\left[ \\alpha_0 - \\alpha_1 C_{N+1} \\lambda_j + \\alpha_2 C_{N+1}^2 \\lambda_j^2 + \\alpha_3 C_{N+1}^2 t_Q \\right], \n\\end{equation}\nwhere\n\\begin{equation}\nj = \\begin{cases}\n1 & C_{N+1} > \\frac{\\alpha_1}{\\alpha_2} \\frac{1}{\\lambda_1 + \\lambda_m} \\text{ or } C_{N+1} \\leq 0 \\\\\nm & \\text{otherwise}.\n\\end{cases}\n\\end{equation}\nMoreover, (\\ref{eqn-sgd-lemma:technical-systematic-component}) is bounded below by\n\\begin{equation}\ne_N\\left[ \\alpha_0 - \\alpha_1 C_{N+1} \\lambda_j + \\alpha_2 C_{N+1}^2 \\lambda_j^2 + \\alpha_3 C_{N+1}^2 s_Q \\right],\n\\end{equation}\nwhere \n\\begin{equation}\nj = \\begin{cases}\n1 & C_{N+1} \\in \\left(0 ,\\frac{\\alpha_1}{\\alpha_2} \\frac{1}{\\lambda_1 + \\lambda_2}\\right] \\\\\nl ~(l \\in \\lbrace 2,\\ldots,m-1 \\rbrace) & C_{N+1} \\in \\left( \\frac{\\alpha_1}{\\alpha_2} \\frac{1}{\\lambda_{l-1} + \\lambda_l}, \\frac{\\alpha_1}{\\alpha_2} \\frac{1}{\\lambda_l + \\lambda_{l+1}} \\right] \\\\\nm & C_{N+1} > \\frac{\\alpha_1}{\\alpha_2} \\frac{1}{\\lambda_{m-1} + \\lambda_m} \\text{ or } C_{N+1} \\leq 0.\n\\end{cases}\n\\end{equation}\n\\end{lemma}\n\\begin{proof}\nWhen $e_N = 0 $, then $e_{N,j} =0$ for $j=2,3$ and $e_{N,M}=0$ by Lemma \\ref{lemma-sgd:higher-moment-bounds}. Therefore, if $e_N = 0$, the bounds hold trivially. So, we will assume that $e_N > 0$. \n\nLet $u_1,\\ldots,u_m$ be the orthonormal eigenvectors corresponding to eigenvalues $\\lambda_1,\\ldots,\\lambda_m$ of $\\E{Q}$. Then, \n\\begin{equation}\n\\frac{e_{N,j}}{e_N} = \\sum_{l=1}^m \\lambda_l^{j-1}\\frac{\\lambda_l (u_l'(\\theta_N - \\theta^*))^2}{e_N}.\n\\end{equation}\nDenote the ratio on the right hand side by $w_l$, and note that $\\lbrace w_l :l=1,\\ldots,m\\rbrace$ sum to one and are nonnegative. With this notation, we bound (\\ref{eqn-sgd-lemma:technical-systematic-component}) from above by\n\\begin{equation}\ne_N\\left[ \\alpha_0 - \\alpha_1 C_{N+1} \\sum_{l=1}^m \\lambda_l w_l + \\alpha_2 C_{N+1}^2 \\sum_{l=1}^m \\lambda_l^2 w_l + \\alpha_3 C_{N+1}^2 t_Q \\right], \n\\end{equation}\nand from below by the same equation but with $t_Q$ replaced by $s_Q$. By the properties of $\\lbrace w_l : l=1,\\ldots, m \\rbrace$, we see that the bounds are composed of convex combinations of quadratics of the eigenvalues. Thus, for the upper bound, if we assign all of the weight to the eigenvalue that maximizes the polynomial $-\\alpha_1 C_{N+1} \\lambda + \\alpha_2 C_{N+1}^2 \\lambda^2$, then we will have the upper bound presented in the result. To compute the lower bounds, we do the analogous calculation.\n\\end{proof}\n\nUsing these two lemmas, we can conclude the following for the homogeneous case.\n\n\\begin{theorem}[Homogeneous Minimizer] \\label{theorem-sgd:homogeneous}\nLet $\\lbrace \\theta_N : N \\in \\mathbb{N} \\rbrace$ be the iterates of SGD-$k$ from Definition \\ref{definition-sgd:sgd-k}. Let $\\theta^*$ be a homogeneous solution to the quadratic problem.\n\nIf $0 < C_{N+1}$ and \n\\begin{equation} \\label{eqn-sgd-theorem:homogeneous-ub}\nC_{N+1} < \\begin{cases}\n\\frac{2\\lambda_1}{\\lambda_1^2 + t_Q\/k} & k > t_Q\/(\\lambda_1 \\lambda_m) \\\\\n\\frac{2\\lambda_m}{\\lambda_m^2 + t_Q\/k} & k \\leq t_Q\/(\\lambda_1 \\lambda_m),\n\\end{cases}\n\\end{equation}\nthen $\\E{e_{N+1}} < \\E{e_N}$. Moreover, if $\\lbrace C_{N+1} \\rbrace$ satisfy (\\ref{eqn-sgd-theorem:homogeneous-ub}) and are uniformly bounded away from zero, then $\\E{e_{N}}$ decays exponentially to zero.\n\nMoreover, there exists a deterministic constant $\\rho_N > 1$ such that $\\cond{e_{N+1}}{\\theta_N} > \\rho_N e_N$ with probability one, if either $C_{N+1} < 0$ or \n\\begin{equation} \\label{eqn-sgd-theorem:homogeneous-lb}\nC_{N+1} > \\frac{2\\lambda_j}{\\lambda_j^2 + s_Q\/k},\n\\end{equation}\nwhere\n\\begin{equation}\nj = \\begin{cases}\n1 & k \\leq \\frac{s_Q}{\\lambda_1\\lambda_2} \\\\\nl ~(l \\in \\lbrace 2,\\ldots,m-1 \\rbrace) & \\frac{s_Q}{\\lambda_l\\lambda_{l-1}} < k \\leq \\frac{s_Q}{\\lambda_{l+1}\\lambda_l} \\\\\nm & k > \\frac{s_Q}{\\lambda_m \\lambda_{m-1}}.\n\\end{cases}\n\\end{equation}\n\\end{theorem}\n\\begin{proof}\nThe upper and lower bounds follow from Lemma \\ref{lemma-sgd:quadratic-recursion} and Lemma \\ref{lemma-sgd:technical-systematic-component}, and observations that $k > t_Q\/(\\lambda_1 \\lambda_m)$ implies\n\\begin{equation} \\label{eqn-sgd-theorem:homogeneous-proof-1}\n\\frac{2\\lambda_j}{\\lambda_j^2 + t_Q\/k} > \\frac{2}{\\lambda_1 + \\lambda_m},\n\\end{equation}\nfor $j = 1,m$. For convergence, when $C_{N+1} > 0$ and strictly less than expression on the left hand side of (\\ref{eqn-sgd-theorem:homogeneous-proof-1}), the upper bound in Lemma \\ref{lemma-sgd:technical-systematic-component} is strictly less than one. Thus, if $\\lbrace C_{N+1} \\rbrace$ are uniformly bounded away from zero, the exponential convergence rate of $\\E{e_N}$ follows trivially.\n\nFor the lower bound, we make note of several facts. First, if $k \\geq \\frac{s_Q}{\\lambda_l \\lambda_{l+1}}$ then $k \\geq \\frac{s_Q}{\\lambda_j \\lambda_{j+1}}$ for $j=1,2,\\ldots,l$. Moreover, when $k \\geq \\frac{s_Q}{\\lambda_l \\lambda_{l+1}}$, then \n\\begin{equation}\n\\frac{2\\lambda_l}{\\lambda_l^2 + s_Q\/k} \\geq \\frac{2}{\\lambda_l + \\lambda_{l+1}}.\n\\end{equation}\nBy Lemma \\ref{lemma-sgd:technical-systematic-component}, the left hand side of this inequality is the lower bound on $C_{N+1}$ that guarantees that $\\E{e_{N+1}} > e_N$ with probability one. Therefore, whenever $k \\geq \\frac{s_Q}{\\lambda_l\\lambda_{l+1}}$, $\\cond{e_N+1}{\\theta_N}$ is not guaranteed to diverge. On the other hand, if $k < \\frac{s_Q}{\\lambda_l \\lambda_{l-1}}$ then $k < \\frac{s_Q}{\\lambda_j \\lambda_{j-1}}$ for $j=l,l+1,\\ldots,m$. Moreover, if $k < \\frac{s_Q}{\\lambda_l \\lambda_{l-1}}$ then\n\\begin{equation}\n\\frac{2\\lambda_l}{\\lambda_l^2 + s_Q\/k} < \\frac{2}{\\lambda_l + \\lambda_{l-1}}.\n\\end{equation}\nTherefore, since the left hand side of this inequality is the lower bound on $C_{N+1}$ that guarantees that there exists a deterministic scalar $ \\rho_N > 1$ such that $\\cond{e_{N+1}}{\\theta_N} > \\rho_N e_N$ with probability one, the lower bound is guaranteed to diverge for $C_{N+1}$ larger than the right hand side.\n\nThus, by the monotonicty of the eigenvalues, for the $l \\in \\lbrace 2,\\ldots,m-1 \\rbrace$ such that $\\frac{s_Q}{\\lambda_l \\lambda_{l-1}} < k \\leq \\frac{s_Q}{\\lambda_l\\lambda_{l+1}}$, $\\cond{e_{N+1}}{\\theta_N} > \\rho_N e_N$ if\n\\begin{equation}\nC_{N+1} > \\frac{2\\lambda_l}{\\lambda_l^2 + s_Q\/k}.\n\\end{equation}\nNote, we can handle the edge cases similarly.\n\\end{proof}\n\\begin{remark}\nBecause $\\rho_N$ is a deterministic scalar quantity, we can conclude that $\\E{e_{N+1}} > \\rho_N \\E{e_N}$. \n\\end{remark}\n\nUsing Theorem \\ref{theorem-sgd:homogeneous}, as $k \\to \\infty$, we recover Theorem \\ref{theorem:gd}. That is, Theorem \\ref{theorem-sgd:homogeneous} contains gradient descent as a special case. Moreover, Theorem \\ref{theorem-sgd:homogeneous} also captures the exponential divergence property that we observed for gradient descent (see Figure \\ref{figure:divergence-gd}). Theorem \\ref{theorem-sgd:homogeneous} also contains additional subtleties. First, it captures SGD-$k$'s distinct phases for convergence and divergence, which depend on the batch size and expected geometry of the quadratic problem as represented by the eigenvalues, $t_Q$ and $s_Q$. Thus, as opposed to classical gradient descent, SGD-$k$ requires a more complex convergence and divergence statement to capture these distinct phases. \n\nWe now state the analogous statement for the inhomogeneous problem. However, we will need the following additional lemmas to bound the behavior of the cross term in Lemma \\ref{lemma-sgd:quadratic-recursion}. Again, our analysis below is more delicate than the standard approach for analyzing SGD. \n\n\\begin{lemma} \\label{lemma-sgd:technical-pseudo-inverse}\nLet $Q$ and $r$ be as in Problem \\ref{problem-sgd:quadratic}. Then, letting $(\\cdot)^\\dagger$ denote the Moore-Penrose pseudo-inverse, \n\\begin{equation}\n\\E{ Q \\E{Q}(Q \\theta^* + r )} = \\E{Q}\\E{Q}^\\dagger\\E{ Q \\E{Q}(Q \\theta^* + r )} .\n\\end{equation}\n\\end{lemma}\n\\begin{proof}\nNote that for any $x \\in \\texttt{row}(\\E{Q})$, $x = \\E{Q}\\E{Q}^\\dagger x$. Thus, we must show that $\\E{Q \\E{Q} (Q \\theta^* + r)}$ is in $\\texttt{row}(\\E{Q})$. Recall, that if there exists a $v \\in \\mathbb{R}^p$ such that $\\E{Q}v = 0$, then $v' \\E{Q} v = 0$. In turn, $v'Qv = 0$ with probability one, which implies that $Q v = 0$ with probability one. Hence, $\\texttt{null}(\\E{Q}) \\subset \\texttt{null}( Q)$ with probability one. Let the set of probability one be denoted by $\\Omega'$. Then, \n\\begin{equation}\n\\texttt{row}( \\E{Q}) = \\texttt{null}(\\E{Q})^\\perp \\supset \\bigcup_{\\omega \\in \\Omega'} \\texttt{null}(Q)^\\perp = \\bigcup_{\\omega \\in \\Omega'} \\texttt{col}( Q) \\supset \\bigcup_{\\omega \\in \\Omega'} \\texttt{col}( Q \\E{Q} ),\n\\end{equation}\nwhere $\\texttt{row}(Q) = \\texttt{col}(Q)$ by symmetry. Therefore, $Q \\E{Q} (Q\\theta^* + r) \\in \\texttt{row}(\\E{Q})$ with probability one. Hence, its expectation is in $\\texttt{row}(\\E{Q})$. \n\\end{proof}\n\n\\begin{lemma} \\label{lemma-sgd:technical-cross-term-bound}\nUnder the setting of Lemma \\ref{lemma-sgd:quadratic-recursion}, for any $\\phi > 0$ and $j \\in \\mathbb{N}$, \n\\begin{equation}\n\\begin{aligned}\n&2 \\left\\vert\\frac{C_{N+1}^2}{k} (\\theta_N - \\theta^*)' \\E{Q \\E{Q}(Q\\theta^* + r)}\\right\\vert \\\\\n&\\quad \\leq 2 \\frac{C_{N+1}^2}{k} \\norm{ \\E{Q}^{j\/2} (\\theta_N - \\theta^*)}_2 \\norm{(\\E{Q}^{j\/2})^\\dagger \\E{Q \\E{Q}(Q\\theta^* + r)}}_2 \\\\\n&\\quad \\leq \\frac{\\phi C_{N+1}^2}{k} (\\theta_N - \\theta^*)' \\E{Q}^j (\\theta_N - \\theta^*) \\\\\n&\\quad \\quad + \\frac{C_{N+1}^2}{\\phi k}\\E{(Q\\theta^* + r)' \\E{Q} Q} ( \\E{Q}^j)^\\dagger \\E{Q \\E{Q}(Q \\theta^* + r)},\n\\end{aligned}\n\\end{equation}\nfor any $N = 0,1,2,\\ldots$.\n\\end{lemma}\n\\begin{proof}\nWe will make use of three facts: Lemma \\ref{lemma-sgd:technical-pseudo-inverse}, H\\\"{o}lder's inequality, and the fact that for any $\\phi > 0$, $ \\varphi \\geq 0$ and $x \\in \\mathbb{R}$, $2 |\\varphi x| \\leq \\phi x^2 + \\phi^{-1} \\varphi^2.$ Using these facts, we have\n\\begin{equation}\n\\begin{aligned}\n&2 \\left\\vert\\frac{C_{N+1}^2}{k} (\\theta_N - \\theta^*)' \\E{Q \\E{Q}(Q\\theta^* + r)}\\right\\vert \\\\\n&\\quad = 2 \\left\\vert\\frac{C_{N+1}^2}{k} (\\theta_N - \\theta^*)'\\E{Q}^{j\/2}(\\E{Q}^{j\/2})^\\dagger \\E{Q \\E{Q}(Q\\theta^* + r)}\\right\\vert\\\\\n&\\quad \\leq 2 \\frac{C_{N+1}^2}{k} \\norm{ \\E{Q}^{j\/2} (\\theta_N - \\theta^*)}_2 \\norm{(\\E{Q}^{j\/2})^\\dagger \\E{Q \\E{Q}(Q\\theta^* + r)}}_2 \\\\\n&\\quad \\leq \\frac{\\phi C_{N+1}^2}{k} (\\theta_N - \\theta^*)' \\E{Q}^j (\\theta_N - \\theta^*) \\\\\n&\\quad + \\frac{C_{N+1}^2}{\\phi k}\\E{(Q\\theta^* + r)' \\E{Q} Q} ( \\E{Q}^j)^\\dagger \\E{Q \\E{Q}(Q \\theta^* + r)}.\n\\end{aligned}\n\\end{equation}\n\\end{proof}\n\n\\begin{lemma} \\label{lemma-sgd:quadratic-recursion-ineq}\nUnder the setting of Lemma \\ref{lemma-sgd:quadratic-recursion},\n\\begin{equation}\n\\begin{aligned}\n\\cond{e_{N+1}}{\\theta_N} &\\leq e_{N} - 2C_{N+1} e_{N,2} + C_{N+1}^2 e_{N,3} + \\frac{C_{N+1}^2}{k} e_{N,M} + \\frac{C_{N+1}^2}{k} e_{N,3} \\\\\n\t\t\t\t\t\t&+ \\frac{C_{N+1}^2}{k} \\E{ (Q\\theta^* + r)' \\E{Q}(Q \\theta^* + r)} \\\\\n\t\t\t\t\t\t&+ \\frac{C_{N+1}^2}{k} \\E{(Q\\theta^*+r)'\\E{Q} Q} ( \\E{Q}^3)^\\dagger \\E{Q \\E{Q}(Q \\theta^* + r)}\n\\end{aligned}\n\\end{equation}\nNow, for any $\\varphi > 0$, $\\exists \\psi \\in \\mathbb{R}$ such that \n\\begin{equation}\n\\begin{aligned}\n\\cond{e_{N+1}}{\\theta_N} &\\geq \\left(1 - \\frac{\\varphi}{k} \\right) e_{N} - 2C_{N+1}e_{N,2} + C_{N+1}^2 e_{N,3} + \\frac{C_{N+1}^2}{k} e_{N,M} \\\\\n\t\t\t\t\t\t&+ \\frac{C_{N+1}^2}{k} \\psi,\n\\end{aligned}\n\\end{equation}\nwhere $\\psi \\geq 0$ if\n\\begin{equation}\n\\varphi\\frac{\\E{ (Q\\theta^* + r)' \\E{Q}(Q \\theta^* + r)}}{\\E{(Q\\theta^*+r)'\\E{Q} Q} \\E{Q}^\\dagger \\E{Q \\E{Q}(Q \\theta^* + r)}} \\geq C_{N+1}^2.\n\\end{equation}\n\\end{lemma}\n\\begin{proof}\nWhen $C_{N+1} = 0,$ the statements hold by Lemma \\ref{lemma-sgd:quadratic-recursion}. Suppose $C_{N+1} \\neq 0$. Using Lemma \\ref{lemma-sgd:quadratic-recursion} and applying Lemma \\ref{lemma-sgd:technical-cross-term-bound} with $\\phi =1$ and $j=3$, the upper bound follows. Now, using Lemma \\ref{lemma-sgd:quadratic-recursion} and applying Lemma \\ref{lemma-sgd:technical-cross-term-bound} with $\\phi = C_{N+1}^{-2} \\varphi$ for $\\varphi > 0$ and $j=1$, we have\n\\begin{equation}\n\\begin{aligned}\n\\cond{e_{N+1}}{\\theta_N} &\\geq \\left(1 - \\frac{\\varphi}{k} \\right) e_{N} - 2C_{N+1} e_{N,2} + C_{N+1}^2 e_{N,3} + \\frac{C_{N+1}^2}{k} e_{N,M} \\\\\n\t\t\t\t\t\t&+ \\frac{C_{N+1}^2}{k} \\psi,\n\\end{aligned}\n\\end{equation}\nwhere\n\\begin{equation}\n\\begin{aligned}\n\\psi &= - \\frac{C_{N+1}^2}{\\varphi} \\E{(Q\\theta^*+r)'\\E{Q} Q}\\E{Q}^\\dagger \\E{Q \\E{Q}(Q \\theta^* + r)} \\\\\n&\\quad + \\E{ (Q\\theta^* + r)' \\E{Q}(Q \\theta^* + r)}.\n\\end{aligned}\n\\end{equation}\nTherefore, if $\\varphi$ is given as selected in the statement of the result, then $\\psi \\geq 0$.\n\\end{proof}\n\n\\begin{theorem}[Inhomogeneous Minimizer] \\label{theorem-sgd:inhomogeneous}\nLet $\\lbrace \\theta_N : N \\in \\mathbb{N} \\rbrace$ be the iterates of SGD-$k$ from Definition \\ref{definition-sgd:sgd-k}. Suppose $\\theta^*$ is an inhomogeneous solution to the quadratic problem.\n\nIf $0 < C_{N+1}$ and\n\\begin{equation} \\label{eqn-theorem-sgd:inhomogeneous-ub}\nC_{N+1} < \\begin{cases}\n\\frac{2\\lambda_1^2}{(1 + 1\/k) \\lambda_1^2 + t_Q\/k } & k+1 > t_Q\/(\\lambda_1 \\lambda_m) \\\\\n\\frac{2\\lambda_m^2}{(1 + 1\/k) \\lambda_m^2 + t_Q\/k } & k+1 \\leq t_Q\/(\\lambda_1\\lambda_m), \n\\end{cases}\n\\end{equation}\nthen $\\E{e_{N+1}} < \\rho_N \\E{e_{N}} + C_{N+1}^2 \\frac{\\psi}{k}$, where $\\rho_N < 1$ and is uniformly bounded away from zero for all $N$, and $\\psi > 0$. Moreover, if $\\lbrace C_{N} : N \\in \\mathbb{N} \\rbrace$ are nonnegative, converge to zero and sum to infinity then $\\E{e_{N}} \\to 0$ as $N \\to \\infty$.\n\nFurthermore, suppose $C_{N+1} < 0$ or \n\\begin{equation} \\label{eqn-theorem-sgd:inhomogeneous-lb}\nC_{N+1} > \\frac{2(\\lambda_j + \\gamma)}{\\lambda_j^2 + s_Q\/k},\n\\end{equation} \nwhere $4\\gamma^2 \\in \\left(0, \\frac{s_Q}{k} \\right]$\nand \n\\begin{equation}\nj = \\begin{cases}\n1 & k \\leq \\frac{s_Q}{ \\lambda_1 \\lambda_2 + \\gamma(\\lambda_1 + \\lambda_2)} \\\\\nl ~(l \\in \\lbrace 2,\\ldots,m-1 \\rbrace) & \\frac{s_Q}{\\lambda_l\\lambda_{l-1} + \\gamma(\\lambda_l + \\lambda_{l-1})} < k \\leq \\frac{s_Q}{\\lambda_l \\lambda_{l+1} + \\gamma (\\lambda_l + \\lambda_{l+1})} \\\\\nm & k > \\frac{s_Q}{\\lambda_m \\lambda_{m-1} + \\gamma(\\lambda_m + \\lambda_{m-1})}.\n\\end{cases}\n\\end{equation}\nThen, $\\exists \\delta > 0$ such that if $\\norm{\\E{Q}(\\theta_{N}-\\theta^*)} < \\delta$ then $\\cond{e_{N+1}}{\\theta_N} \\geq \\rho_N e_{N} + C_{N+1}^2 \\psi$ with probability one, where $\\rho_N > 1$ and $\\psi > 0$ independently of $N$.\n\\end{theorem}\n\\begin{proof}\nFor the upper bound, we apply Lemmas \\ref{lemma-sgd:quadratic-recursion-ineq} and \\ref{lemma-sgd:technical-systematic-component} to conclude that\n\\begin{equation}\n\\cond{e_{N+1}}{\\theta_N} \\leq e_{N}\\left(1 - 2 C_{N+1}\\lambda_j + \\left(1 + \\frac{1}{k}\\right) C_{N+1}^2\\lambda_j^2 + C_{N+1}^2 \\frac{t_Q}{k} \\right) + C_{N+1}^2 \\frac{\\psi}{k},\n\\end{equation}\nwhere $\\psi > 0$ is independent of $N$ and\n\\begin{equation}\nj = \\begin{cases}\n1 & C_{N+1} \\leq 0 \\text{ or } \\frac{2}{(1 + 1\/k)(\\lambda_1 + \\lambda_m)} \\\\\nm & \\text{otherwise}.\n\\end{cases}\n\\end{equation}\nMoreover, the component multiplying $e_N$ is less than one if $C_{N+1} > 0$ and\n\\begin{equation} \\label{eqn-sgd-theorem:proof-inhomo-1}\nC_{N+1} < \\frac{2\\lambda_j}{\\left( 1 + \\frac{1}{k}\\right)\\lambda_j^2 + \\frac{t_Q}{k}}.\n\\end{equation}\nWhen $k + 1 > t_Q\/(\\lambda_1\\lambda_m)$, the right hand side of (\\ref{eqn-sgd-theorem:proof-inhomo-1}) is larger than $\\frac{2}{(1 + 1\/k)(\\lambda_1 + \\lambda_m)}$. The upper bound follows. Convergence follows by applying Lemma 1.5.2 of \\cite{bertsekas1999}.\n\nFor the lower bound, note that since $\\theta^*$ is inhomogeneous, $\\Prb{Q = \\E{Q}} < 1$. Therefore, $s_Q > 0$. Now, we apply Lemmas \\ref{lemma-sgd:quadratic-recursion-ineq} and \\ref{lemma-sgd:technical-systematic-component} to conclude that for any $\\gamma > 0$ there exists $\\psi_N$ such that\n\\begin{equation}\n\\cond{e_{N+1}}{\\theta_N} \\geq e_{N}\\left(\\left(1 - \\frac{4\\gamma^2}{\\lambda_j^2 + \\frac{1}{k}s_Q} \\right) - 2 C_{N+1}\\lambda_j + C_{N+1}^2\\lambda_j^2 + C_{N+1}^2 \\frac{s_Q}{k} \\right) + C_{N+1}^2 \\frac{\\psi_N}{k},\n\\end{equation}\nwhere $\\psi_N$ are nonnegative and uniformly bounded from zero for $C_{N+1}$ sufficiently small, and\n\\begin{equation}\nj = \\begin{cases}\n1 & C_{N+1} \\in \\left(0, \\frac{2}{\\lambda_1 + \\lambda_2} \\right] \\\\\nl ~(l \\in \\lbrace 2,\\ldots,m-1 \\rbrace) & C_{N+1} \\in \\left( \\frac{2}{\\lambda_{l-1} + \\lambda_l}, \\frac{2}{\\lambda_l + \\lambda_{l+1}} \\right] \\\\\nm & C_{N+1} \\leq 0 \\text{ or } \\frac{2}{\\lambda_m + \\lambda_{m-1}} < C_{N+1}.\n\\end{cases}\n\\end{equation}\nThe term multiplying $e_N$ can be rewritten as\n\\begin{equation} \\label{eqn-sgd-theorem:proof-inhomo-2}\n1 - \\frac{4\\gamma^2}{\\lambda_j^2 + \\frac{1}{k}s_Q} - \\frac{\\lambda_j^2}{\\lambda_j^2 + \\frac{1}{k}s_Q} + \\left(\\lambda_j^2 + \\frac{1}{k}s_Q\\right)\\left( C_{N+1} - \\frac{\\lambda_j}{\\lambda_j^2 + \\frac{1}{k}s_Q} \\right)^2,\n\\end{equation}\nwhich nonnegative when $4\\gamma^2$ is in the interval given by in the statement of the result. Moreover, if\n\\begin{equation} \\label{eqn-sgd-theorem:proof-inhomo-3}\nC_{N+1} > \\frac{2(\\lambda_j + \\gamma)}{\\lambda_j^2 + \\frac{s_Q}{k}},\n\\end{equation}\nthen\n\\begin{equation}\nC_{N+1}\t> \\frac{\\lambda_j}{\\lambda_j^2 + \\frac{s_Q}{k}} + \\frac{\\sqrt{\\lambda_j^2 + 4\\gamma^2}}{\\lambda_j^2 + \\frac{s_Q}{k}},\t\n\\end{equation}\nwhich implies\n\\begin{equation}\n\\left(C_{N+1} - \\frac{\\lambda_j}{\\lambda_j^2 + \\frac{s_Q}{k}} \\right)^2 > \\frac{\\lambda_j^2 + 4\\gamma^2}{\\left(\\lambda_j^2 + \\frac{s_Q}{k}\\right)^2}.\n\\end{equation}\nFrom this inequality, we conclude that if (\\ref{eqn-sgd-theorem:proof-inhomo-3}) holds then (\\ref{eqn-sgd-theorem:proof-inhomo-2}) is strictly larger than $1$. Lastly, for any $\\tilde{j} \\neq j$, \n\\begin{equation}\n\\frac{s_Q}{\\lambda_j \\lambda_{\\tilde{j}} + \\gamma(\\lambda_j + \\lambda_{\\tilde{j}})} > k\n\\end{equation}\nis equivalent to\n\\begin{equation}\n\\frac{2}{\\lambda_j + \\lambda_{\\tilde{j}}} > \\frac{2(\\lambda_j + \\gamma)}{\\lambda_j^2 + \\frac{1}{k}s_Q}.\n\\end{equation}\nWith these facts, the conclusion for the lower bound follows just as in the proof of Theorem \\ref{theorem-sgd:homogeneous}.\n\\end{proof}\n\nComparing Theorems \\ref{theorem-sgd:homogeneous} and \\ref{theorem-sgd:inhomogeneous}, we see that inhomogeneity exacts an additional price. To be specific, call the term multiplying $e_N$ in the bounds as the systematic component and the additive term as the variance component. First, the inhomogeneous case contains a variance component that was not present for the homogeneous case. Of course, this variance component is induced by the inhomogeneity: at each update, SGD-$k$ is minimizing an approximation to the quadratic problem that has a distinct solution from $\\theta^*$ with nonzero probability, which requires decaying the step sizes to remove the variability of solutions induced by these approximations. Second, the inhomogeneity shrinks the threshold for step sizes that ensure a reduction in the systematic component, and inflates the threshold for step sizes that ensure an increase in the systematic component.\n\nContinuing to compare the results for the homogeneous and inhomogeneous cases, we also observe that the systematic component diverges exponentially when the expected gradient is sufficiently small (less than $\\delta$). Indeed, this is the precise region in which such a result is important: it says that for step sizes that are too large, any estimate near the solution set is highly unstable. In other words, if we initialize SGD-$k$ near the solution set, our iterates will diverge exponentially fast from this region for $C_N$ above the divergence threshold.\n\nTo summarize, our analysis of the stochastic quadratic problem defined in Problem \\ref{problem-sgd:quadratic} is a generalization of the standard quadratic problem considered in classical numerical optimization. Moreover, our analysis of SGD-$k$ (see Definition \\ref{definition-sgd:sgd-k}) on the stochastic quadratic problem generalizes the results and properties observed for gradient descent. In particular, our analysis (see Theorems \\ref{theorem-sgd:homogeneous} and \\ref{theorem-sgd:inhomogeneous}) derive a threshold for step sizes above which the iterates will diverge, and this threshold includes the threshold for classical gradient descent. Moreover, our analysis predicts an exponential divergence of the iterates from the minimizer. \n\n\\subsection{The Deterministic Mechanism} \\label{subsection:deterministic-mechanism}\n\nWe now use the insight developed from our analysis of the generic quadratic problem to formulate a deterministic mechanism for how SGD-$k$ ``escapes'' local minimizers. We will begin by stating the generic nonconvex problem, formalizing our deterministic mechanism, and discussing its implications for nonconvex optimization. We then briefly discuss the pitfalls of naive attempts to directly apply the theory developed in \\S \\ref{subsection:analysis} to arbitrary nonconvex problems.\n\n\\begin{problem}[Nonconvex Problem] \\label{problem:nonconvex}\nLet $(\\Omega,\\mathcal{F},\\mathbb{P})$ be a probability space, and \nlet $C^2(\\mathbb{R}^p,\\mathbb{R})$ be the set of twice continuously differentiable functions from $\\mathbb{R}^p$ to\nMoreover, let $f$ be a random quantity taking values in $C^2(\\mathbb{R}^p,\\mathbb{R})$. Moreover, suppose that\n\\begin{enumerate}\n\\item $\\E{f(\\theta)}$ exists for all $\\theta \\in \\mathbb{R}^p$, and is denoted by $F(\\theta)$.\n\\item Suppose $F$ is bounded from below. \n\\item Suppose that $F \\in C^2(\\mathbb{R}^p,\\mathbb{R})$, and $\\nabla F(\\theta) = \\E{ \\nabla f(\\theta) }$ and $\\nabla^2 F(\\theta) = \\E{ \\nabla^2 f(\\theta)}$. \n\\end{enumerate}\nLet $\\lbrace f_N: N \\in \\mathbb{N} \\rbrace$ be independent copies of $f$. The nonconvex problem is to use $\\lbrace f_N : N \\in \\mathbb{N} \\rbrace$ to determine a $\\theta^*$ that is a minimizer of $F(\\theta)$. \n\\end{problem}\n\nNote, while Problem \\ref{problem:nonconvex} focuses on twice differentiable functions, we often only have continuous functions in machine learning applications. However, for most machine learning applications, the continuous functions are only nondifferentiable at finitely many points, which can be approximated using twice continuously differentiable functions.\n\nBefore stating the deterministic mechanism, we will also need to extend our definition of SGD-$k$ to the generic nonconvex problem. Of course, this only requires replacing (\\ref{eqn-definition:sgd-k}) with \n\\begin{equation}\n\\theta_{N+1} = \\theta_N - \\frac{C_{N+1}}{k}\\sum_{j=Nk+1}^{N(k+1)} \\nabla f_j(\\theta_N).\n\\end{equation}\nWe now state our deterministic mechanism. In this statement, $B(\\theta,\\epsilon)$ denotes a ball around a point $\\theta \\in \\mathbb{R}^p$ of radius $\\epsilon > 0$.\n\n\\begin{definition}[Deterministic Mechanism] \\label{definition:deterministic mechanism}\nLet $\\lbrace \\theta_N : N \\in \\mathbb{N} \\rbrace$ be the iterates generated by SGD-$k$ with deterministic step sizes $\\lbrace C_N : N \\in \\mathbb{N} \\rbrace$ for the nonconvex problem defined in Problem \\ref{problem:nonconvex}. Suppose $\\theta^*$ is a local minimizer of Problem \\ref{problem:nonconvex}.\n\nFor an $\\epsilon > 0$ sufficiently small, let $\\lambda_1 \\geq \\cdots \\geq \\lambda_m > 0$ denote the nonzero eigenvalues of\n\\begin{equation} \n\\frac{1}{|B(\\theta^*,\\epsilon)|} \\int_{B(\\theta^*,\\epsilon)}\\nabla^2 F(\\theta) d\\theta,\n\\end{equation}\nand define $s_f$ to be the maximum between $0$ and \n\\begin{equation}\n\\frac{1}{|B(\\theta^*,\\epsilon)|}\\int_{B(\\theta^*,\\epsilon)}\\inf_{v'\\nabla^2 F(\\theta) v \\neq 0\n} \\left\\lbrace \\frac{v' \\left(\\E{ \\nabla^2 f(\\theta) \\nabla^2 F(\\theta) \\nabla^2 f(\\theta)} - (\\nabla^2 F(\\theta) )^3\\right) v}{v' \\nabla^2 F(\\theta) v} \\right\\rbrace d\\theta.\n\\end{equation}\nThere is a $\\delta > 0$ such that if $\\theta_N \\in B(\\theta^*,\\delta) \\setminus \\lbrace \\theta^* \\rbrace$ and $C_{N+1} < 0$ or\n\\begin{equation}\nC_{N+1} > \\frac{2\\lambda_j}{\\lambda_j^2 + \\frac{1}{k}s_f},\n\\end{equation} \nwhere\n\\begin{equation}\nj = \\begin{cases}\n1 & k \\leq \\frac{s_f}{ \\lambda_1 \\lambda_2 } \\\\\nl ~(l \\in \\lbrace 2,\\ldots,m-1 \\rbrace) & \\frac{s_f}{\\lambda_l\\lambda_{l-1}} < k \\leq \\frac{s_f}{\\lambda_l \\lambda_{l+1}} \\\\\nm & k > \\frac{s_f}{\\lambda_m \\lambda_{m-1}},\n\\end{cases}\n\\end{equation}\nthen for some $\\rho_N > 1$, $\\cond{F(\\theta_{N+1}) - F(\\theta^*)}{\\theta_N} > \\rho_N ( F(\\theta_N) - F(\\theta^*) )$ with probability one.\n\\end{definition}\n\nFor nonconvex problems, our deterministic mechanism makes several predictions. First, suppose a minimizer's local geometry has several large eigenvalues. Then, for $k$ sufficiently small, SGD-$k$'s divergence from a local minimizer will be determined by these large eigenvalues. This prediction is indeed observed in practice: \\cite{keskar2016} noted that methods that were ``more stochastic'' (i.e., small $k$) diverged from regions that had only a fraction of relatively large eigenvalues. \n\nSecond, suppose for a given nonconvex function, we have finitely many local minimizers and, for each $k$, we are able to compute the threshold for each local minimizer (for some sufficiently small $\\epsilon > 0$, which may differ for each minimizer). Then, for identical choices in step sizes, $\\lbrace C_{N+1} : N \\in \\mathbb{N} \\rbrace$, the number of minimizers for which SGD-$k$ can converge to is a nondecreasing function of $k$. Moreover, SGD-$k$ and Gradient Descent will only converge to minimizers that can support the step sizes $\\lbrace C_{N+1} : N \\in \\mathbb{N} \\rbrace$, which implies that SGD-$k$ will prefer flatter minimizers in comparison to gradient descent. Again, this prediction is indeed observed in practice: the main conclusion of \\cite{keskar2016} is that SGD-$k$ with smaller $k$ converges to flatter minimizers in comparison to SGD-$k$ with large $k$. \n\nThird, our deterministic mechanism states that whenever the iterates are within some $\\delta$ ball of $\\theta^*$ with step sizes that are too large, we will observe exponential divergence. Demonstrating that exponential divergence does indeed occur in practice is the subject of \\S \\ref{section:homogeneous} and \\S \\ref{section:inhomogeneous}.\n\nHowever, before discussing these experiments, it is worth noting the challenges of generalizing the quadratic analysis to the nonconvex case. In terms of convergence, recall that a key ingredient in the local quadratic analysis for classical numerical optimization is stability: once iterates enter a basin of the minimizer, then the iterates will remain within this particular basin of the minimizer. However, for stochastic problems, this will not hold with probability one, especially when the gradient has a noise component with an unbounded support (e.g., linear regression with normally distributed errors). \n\nFortunately, attempts to generalize the divergence results will not require this stability property because divergence implies that the iterates will eventually escape from any such basin. The challenge for divergence, however, follows from the fact that the lower bounds are highly delicate, and any quadratic approximation will require somewhat unrealistic conditions on the remainder term in Taylor's theorem for the nonconvex problem. Thus, generalizing the quadratic case to the nonconvex case is something that we are actively pursuing. \n\\subsection{The Stochastic Mechanism}\n\nConsider using SGD on a one-dimensional stochastic nonconvex problem whose expected objective function is given by the graph in Figure \\ref{figure:one-dim-obj-A}. The stochasticity of the nonconvex problem is generated by adding independent symmetric, bounded random variables to the derivative evaluation at any point. That is, at any iteration, SGD's search direction is generated by the derivative of the expected objective function plus some symmetric, bounded random noise. \n\nIntuitively, the noise in the SGD search directions, by the size of the basins alone, will force the iterates to spend more time in the flatter minimizer on the left. That is, the noise of in SGD's search directions will increase an iterates probability of being in the flatter basin in comparison to the sharper basin. This is precisely stochastic mechanism.\n\n\\begin{figure}[bh]\n\\centering\n\\input{figure\/sgd_stochastic_mechanism_A}\n\\caption{The expected objective function for a simple one-dimensional nonconvex problem. The minimizer on the left is the flatter minimizer. The minimizer on the right is the sharper minimizer.}\n\\label{figure:one-dim-obj-A}\n\\end{figure}\n\nFor another example, consider a similar one-dimensional stochastic nonconvex problem whose expected objective function is given by the graph in Figure \\ref{figure:one-dim-obj-B}. In this example, the minimizer on the left is the flatter minimizer while the minimizer on the right is the sharper minimizer. Now the stochastic mechanism would predict that the probability of an iterate being in the basin of the sharper minimizer is actually greater than that of the flatter minimizer. \n\n\\begin{figure}\n\\centering\n\\input{figure\/sgd_stochastic_mechanism_B}\n\\caption{The expected objective function for a different one-dimensional nonconvex problem. The minimizer on the left is the flatter minimizer. The minimizer on the right is the sharper minimizer.}\n\\label{figure:one-dim-obj-B}\n\\end{figure} \n\nThus, the stochastic mechanism is not explaining why SGD prefers flatter minimizers, but rather it is explaining why SGD has a higher probability of being in basins with larger size. Moreover, while flatness and basin size might seem to be intuitively connected, as Figure \\ref{figure:one-dim-obj-B} shows, this is not necessarily the case. Therefore, according to this example, we might assume that flatness is an inappropriate term, and it is actually the volume of the minimizer that is important (based on the stochastic mechanism) in determining to which minimizers the SGD iterates converge.\n\nHowever, ignoring flatness or sharpness of a minimizer in favor of the size of its basin of attraction does not align with another experimental observation: in high-dimensional problems, SGD even avoids those sharper minimizers that had a small fraction of their eigenvalues termed relatively ``large'' \\cite{keskar2016}. Importantly, in such high-dimensions, the probability of the noise exploring those limited directions that correspond to the ``larger'' eigenvalues is low (e.g., consider isotropic noise in large dimensions). This suggests that, while the stochastic mechanism is intuitive and can explain why SGD may prefer large basins of attractions, it is insufficient to explain why SGD prefers flatter minimizers over sharper minimizers. Accordingly, what is an appropriate mechanism to describe how SGD iterates ``select'' minimizers?\n\n\\subsection{Main Contribution}\nIn this work, we derive and propose a deterministic mechanism based on the expected local geometry of a minimizer and the batch size of SGD to describe why SGD iterates will converge or diverge from the minimizer. As a consequence of our deterministic mechanism, we can explain why SGD prefers flatter minimizers to sharper minimizers in comparison to gradient descent, and we can explain the manner in which SGD will ``escape'' from certain minimizers. Importantly, we verify the predictions of our deterministic mechanism with numerical experiments on two nontrivial nonconvex problems. \n\n\\subsection{Outline}\n\nIn \\S \\ref{section:classical}, we begin by reviewing the classical mechanism by which gradient descent diverges from a quadratic basin. In \\S \\ref{section:stochastic-quadratic}, we will rigorously generalize the classical, deterministic mechanism to SGD-$k$ (i.e., SGD with a batch size of $k$) on a generic stochastic quadratic problem, and, using this rigorous analysis, we define our deterministic mechanism. In \\S \\ref{section:homogeneous} and \\S \\ref{section:inhomogeneous}, we numerically verify the predictions of our deterministic mechanism on two nontrivial, nonconvex problem. In \\S \\ref{section:conclusion}, we conclude this work. \n\n\n\\section{Introduction}\n\\input{section\/introduction}\n\n\\section{Classical Gradient Descent} \\label{section:classical}\n\\input{section\/classical_gd}\n\n\\section{A Deterministic Mechanism} \\label{section:stochastic-quadratic}\n\\input{section\/quadratic}\n\n\\section{Numerical Study of a Homogeneous Nonconvex Problem} \\label{section:homogeneous}\n\\input{section\/homogeneous}\n\n\\section{Numerical Study of an Inhomogeneous Nonconvex Problem} \\label{section:inhomogeneous}\n\\input{section\/inhomogeneous}\n\n\\section{Conclusion} \\label{section:conclusion}\n\\input{section\/conclusion}\n\n\\section*{Acknowledgements}\nWe would like to thank Rob Webber for pointing out the simple proof item 2 in Lemma \\ref{lemma-sgd:higher-moment-bounds}. We would also like to thank Mihai Anitescu for his general guidance throughout the preparation of this work.\n\n\\section*{Funding}\nThe author is supported by the NSF Research and Training Grant \\# 1547396.\n\n\\bibliographystyle{tfs}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Conclusion}\nIn this paper, we presented MasterGNN, a joint air quality and weather prediction model that explicitly models the correlations and interactions between two predictive tasks.\nSpecifically, we first proposed a heterogeneous recurrent graph neural network to capture spatiotemporal autocorrelation among air quality and weather monitoring stations. Then, we developed a multi-adversarial learning framework to resist observation noise propagation. In addition, an adaptive training strategy is devised to automatically balance the optimization of multiple discriminative losses in multi-adversarial learning. Extensive experimental results on two real-world datasets demonstrate that the performance of MasterGNN on both air quality and weather prediction tasks consistently outperforms seven state-of-the-art baselines.\n\\section{Experiments}\n\n\\begin{table}[t]\n\\centering\n\\begin{small}\n\\caption{Statistics of datasets.}\n\\label{table-data-stats}\n\\vspace{-6pt}\n\\begin{tabular}{c | c c} \\hline\n{Data description} & Beijing & Shanghai \\\\ \\hline \\hline\n\\# of air quality stations & 35 & 76 \\\\ \n\\hline\n\\# of weather stations & 18 & 11\\\\ \n\\hline\n\\# of air quality observations & 409,853 & 1,331,520 \\\\ \n\\hline\n\\# of weather observations & 210,824 & 192,720\\\\ \n\\hline\n\\# of POIs & 900,669 & 1,061,399\\\\ \n\\hline\n\\# of road segments & 812,195 & 768,336 \\\\ \n\\hline\n\\end{tabular}\n\\end{small}\n\\vspace{-7mm}\n\\end{table}\n\n\\begin{table*}[t]\n\\centering\n\\vspace{-10pt}\n\\begin{small}\n\\caption{Overall performance comparison of different approaches. A smaller value indicates a better performance.}\n\\vspace{-5pt}\n\\label{tab-eff}\n\\begin{tabular}{c|c|c|c|c|c|c|c|c}%\n\\hline\n\\multirow{3}{*}{Model} & \\multicolumn{8}{|c}{MAE\/SMAPE}\\\\\n\\cline{2-9}\n\\multirow{3}{*}{} & \\multicolumn{4}{|c}{Beijing} & \\multicolumn{4}{|c}{Shanghai}\\\\\n\\cline{2-9}\n\\multirow{3}{*}{} & {AQI} & {Temperature} & {Humidity} & {Wind speed} & {AQI} & {Temperature} & {Humidity} & {Wind speed}\\\\\n\\hline\nARIMA & 54.97\/0.925 & 3.26\/0.475 & 16.42\/0.403 & 1.17\/0.637 & 34.36\/0.597 & 2.27\/0.394 & 11.96\/0.267 & 1.35\/0.653 \\\\\nLR & 45.27\/0.732 & 2.39\/0.426 & 9.58\/0.296 & 1.04\/0.543 & 27.68\/0.462 & 1.93\/0.342 & 9.53\/0.254 & 1.21\/0.568 \\\\\nGBRT & 38.36\/0.678 & 2.31\/0.397 & 8.43\/0.254 & 0.93\/0.489 & 24.13\/0.423 & 1.87\/0.329 & 8.91\/0.232 & 1.05\/0.497 \\\\\n\\hline\nDeepAir & 31.45\/0.613 & 2.21\/0.375 & 8.25\/0.246 & 0.82\/0.468 & 18.79\/0.297 & 1.79\/0.304 & 7.82\/0.209 & 0.93\/0.516 \\\\\nGeoMAN & 30.04\/0.595 & 2.05\/0.356 & 8.01\/0.223 & 0.71\/0.457 & 20.75\/0.348 & 1.56\/0.261 & 6.45\/0.153 & 0.82\/0.465 \\\\\nGC-DCRNN & 29.58\/0.586 & 2.14\/0.364 & 8.14\/0.235 & 0.73\/0.462 & 18.73\/0.305 & 1.61\/0.263 & 6.83\/0.186 & 0.78\/0.462 \\\\\nDUQ & 32.69\/0.624 & 2.02\/0.348 & 7.97\/0.217 & 0.69\/0.446 & 22.81\/0.364 & 1.42\/0.235 & 6.23\/0.136 & 0.67\/0.442 \\\\\n\\hline\\hline\nMasterGNN & \\textbf{27.45\/0.548} & \\textbf{1.87\/0.326} & \\textbf{7.25\/0.184}& \\textbf{0.64\/0.428} &\\textbf{16.51\/0.265} & \\textbf{1.25\/0.213} & \\textbf{5.64\/0.095}&\\textbf{0.61\/0.427}\\\\\n\\hline\n\\end{tabular}\n\\vspace{-8pt}\n\\end{small}\n\\end{table*}\n\n\\subsection{Experimental settings}\n\\subsubsection{Data description.}\nWe conduct experiments on two real-world datasets collected from Beijing and Shanghai, two metropolises in China. \nThe Beijing dataset is ranged from January 1, 2017, to April 1, 2018, and the Shanghai dataset is ranged from June 1, 2018, to June 1, 2020.\nBoth datasets include (1) air quality observations~(\\emph{i.e.},\\xspace PM2.5, PM10, O3, NO2, SO2, and CO), (2) weather observations and weather forecast~(\\emph{i.e.},\\xspace weather condition, temperature, pressure, humidity, wind speed and wind direction), and (3) urban contextual data~(\\emph{i.e.},\\xspace POI distribution and Road network distribution~\\cite{zhu2016days,liu2019hydra,liu2020polestar}). \nThe observation data of Beijing dataset is from KDD CUP 2018\\footnote{https:\/\/www.biendata.xyz\/competition\/kdd\\_2018\/data\/}, and the observation data of Shanghai dataset is crawled from China government websites\\footnote{http:\/\/www.cnemc.cn\/en\/}.\nAll air quality and weather observations are collected by every hour. \nWe associate POI and road network features to each station through an open API provided by an anonymized online map application.\nSame as existing studies~\\cite{yi2018deep,wang2019deep}, we focus on Air Quality Index~(AQI) for air quality prediction, which is derived by the Chinese AQI standard, and temperature, humidity and wind speed for weather prediction.\nWe split the data into training, validation, and testing data by the ratio of 8:1:1.\nThe statistics of two datasets are shown in Table~\\ref{tab-eff}.\n\n\\subsubsection{Implementation details.}\nOur model and all the deep learning baselines are implemented with PyTorch. All methods are evaluated on a Linux server with 8 NVIDIA Tesla P40 GPUs.\nWe set context-aware heterogeneous graph attention layers $l=2$.\nThe cell size of the GRU is set to $64$. \nWe set $\\epsilon=15$, and the hidden size of the multi-layer perceptron of each discriminator is fixed to $64$. All the learning parameters are initialized with a uniform distribution, and the model is trained by stochastic gradient descent with a learning rate $lr=0.00001$. The activation function used in the hidden layers is LeakyReLU ($\\alpha$=0.2).\nAll the baselines use the same input features as ours, except ARIMA excludes contextual features. \nEach numerical feature is normalized by Z-score. We set $T=72$ and the future time step $\\tau=48$.\nFor a fair comparison, we carefully fine-tuned the hyper-parameters of each baseline, and use one another historical observations as a part of feature input.\n\n\n\\subsubsection{Evaluation metrics.}\nWe use Mean Absolute Error (MAE) and Symmetric Mean Absolute Percentage Error (SMAPE), two widely used metrics~\\cite{luo2019accuair} for evaluation.\n\n\\subsubsection{Baselines.}\nWe compare the performance of MasterGNN with the following seven baselines:\n\n\\begin{itemize}\n \\item \\textbf{ARIMA} ~\\cite{brockwell2002introduction} is a classic time series prediction method utilizes the moving average and auto-regressive component to estimate future trends.\n \\item \\textbf{LR} uses linear regression ~\\cite{montgomery2012introduction} for joint prediction. We concatenate previous observations and contextual features as input.\n \n \\item \\textbf{GBRT} is a tree-based model widely used for regression tasks. We adopt an efficient version XGBoost~\\cite{chen2016xgboost}, and the input feature is same as LR.\n \\item \\textbf{GeoMAN}~\\cite{liang2018geoman} is a multi-level attention-based network which integrates spatial and temporal attention for geo-sensory time series prediction.\n \\item \\textbf{DeepAir}~\\cite{yi2018deep} is a deep learning method for air quality prediction. It combines spatial transformation and distributed fusion components to fuse various urban data.\n \\item \\textbf{GC-DCRNN}~\\cite{lin2018exploiting} leverages a diffusion convolution recurrent neural network to model spatiotemporal dependencies for air quality prediction.\n \\item \\textbf{DUQ}~\\cite{wang2019deep} is an advanced weather forecasting framework that improves numerical weather prediction by using uncertainty quantification.\n\\end{itemize}\n\n\\begin{figure*}[t]\n\t\\begin{minipage}{1.\\linewidth}\n\t\t\\centering\n\t\t\\vspace{0ex}\n\t\n\t\t\\subfigure[{\\small Effect of joint prediction.}]{\\label{copred}\n\t\t\t\\includegraphics[width=0.24\\textwidth]{exp\/exp_copred.pdf}}\n\t\t\\subfigure[{\\small Effect of heterogeneous recurrent graph neural network.}]{\\label{heterst}\n\t\t\t\\includegraphics[width=0.24\\textwidth]{exp\/exp_st.pdf}}\n\t\t\\subfigure[{\\small Effect of multi-adversarial learning.}]{\\label{adv}\n\t\t\t\\includegraphics[width=0.24\\textwidth]{exp\/exp_adv.pdf}}\n\t\t\\subfigure[{\\small Effect of multi-task adaptive training.}]{\\label{opti}\n\t\t\t\\includegraphics[width=0.24\\textwidth]{exp\/exp_ms.pdf}}\n\t\\end{minipage}\n\t\\vspace{-3ex}\n\t\\caption{Ablation study of MasterGNN on the Beijing dataset.} \n\t\\vspace{-7mm}\n\t\\label{parasensitivity}\n\\end{figure*}\n\n\\subsection{Overall performance}\nTable 2 reports the overall performance of MasterGNN and compared baselines on two datasets with respect to MAE and SMAPE. \nWe observe our approach consistently outperforms all the baselines using both metrics, which shows the superiority of MasterGNN.\nSpecifically, our model achieves (7.8\\%, 8.0\\%, 10.0\\%, 7.8\\%) and (6.9\\%, 6.7\\%, 17.9\\%, 4.2\\%) improvements beyond the best baseline on MAE and SMAPE on Beijing for (AQI, temperature, humidity, wind speed) prediction, respectively. Similarly, the improvement of MAE and SMAPE on Shanghai are (13.4\\%, 13.6\\%, 10.4\\%, 9.8\\%) and (12.1\\%, 10.3\\%, 43.2\\%, 3.5\\%), respectively. \nMoreover, we observe all deep learning based models outperform statistical learning based approaches by a large margin, which demonstrate the effectiveness of deep neural network for modeling spatiotemporal dependencies. \nRemarkably, GC-DCRNN performs better than all other baselines for air quality prediction, and DUQ consistently outperforms other baselines for weather prediction. \nHowever, they perform relatively poorly on the other task, which indicates their poor generalization capability and further demonstrate the advantage of joint prediction.\n\n\\begin{figure*}[t]\n\t\\begin{minipage}{1.\\linewidth}\n\t\t\\centering\n\t\t\\vspace{-10pt}\n\t\n\t\t\\subfigure[{\\small Effect of $T$.}]{\\label{length}\n\t\t\t\\includegraphics[width=0.24\\textwidth]{exp\/para_t.pdf}}\n\t\t\\subfigure[{\\small Effect of $d$.}]{\\label{hidden_units}\n\t\t\t\\includegraphics[width=0.24\\textwidth]{exp\/para_d.pdf}}\n\t\t\\subfigure[{\\small Effect of $\\epsilon$.}]{\\label{epsilon}\n\t\t\t\\includegraphics[width=0.24\\textwidth]{exp\/para_epsilon.pdf}}\n\t\t\\subfigure[{\\small Effect of $\\tau$.}]{\\label{tau}\n\t\t\t\\includegraphics[width=0.24\\textwidth]{exp\/para_tau.pdf}}\n\t\\end{minipage}\n\t\\vspace{-3ex}\n\t\\caption{Parameter sensitivities on the Beijing dataset.} \n\t\\vspace{-3ex}\n\t\\label{parasensitivity}\n\\end{figure*}\n\n\\subsection{Ablation study}\nThen we study the effectiveness of each module in MasterGNN. Due to the page limit, we only report the results on Beijing by using SMAPE. The results on Beijing using MAE and on Shanghai using both metrics are similar.\n\n\\textbf{Effect of joint prediction}. To validate the effectiveness of joint prediction, we compare the following variants of MasterGNN:\n(1) SLE models two tasks independently without information sharing,\n(2) CXT models two tasks independently but uses the other observations as a part of contextual features. \nAs depicted in Figure \\ref{copred}, we observe a performance gain of CXT by adding the other observation features, but MasterGNN achieves the best performance.\nSuch results demonstrate the advantage of modeling the interactions between two tasks compared with simply adding the other observations as the feature input.\n\n\\textbf{Effect of heterogeneous recurrent graph neural network}. We evaluate three variants of HRGNN, (1) RS removes the spatial autocorrelation block from MasterGNN, (2) AP replaces GRU by average pooling, (3) RH handles air quality and weather observations homogeneously. As shown in Figure \\ref{heterst}, we observe MasterGNN achieves the best performance compared with other variants, demonstrating the benefits of incorporating spatial and temporal autocorrelations for joint air quality and weather prediction. \n\n\\textbf{Effect of multi-adversarial learning}.\nWe compare the following variants: (1) STD removes the macro discriminator, (2) SAD removes the micro temporal discriminator, (3) TAD removes the micro spatial discriminators. As shown in Figure \\ref{adv}, the macro discriminator is the most influential one, which indicates the importance of defending noises from the global view.\nOverall, MasterGNN consistently outperforms STD, SAD, and TAD on all tasks by integrating all discriminators, demonstrating the benefits of developing multiple discriminators for joint prediction. \n\n\\textbf{Effect of multi-task adaptive training}.\nTo verify the effect of multi-task adaptive training, we develop (1) Average Loss Minimization (ALM) uses the same weight for all discriminators, (2) Fixed Loss Minimization~(FLM) sets weighted but fixed importance based on the performance loss of MasterGNN when the corresponding discriminator is removed.\nAs shown in Figure \\ref{opti}, we observe performance degradation when we fix each discriminator's weight, either on average or weighted. \nOverall, adaptively re-weighting each discriminator's importance can guide the optimization direction of the generator and lead to better performance. \n\n\n\\subsection{Parameter sensitivity}\n\nWe further study the parameter sensitivity of MasterGNN.\nWe report SMAPE on the Beijing dataset. Each time we vary a parameter, we set others to their default values.\n\nFirst, we vary the input length $T$ from 6 to 96. The results are reported in Figure ~\\ref{length}. As the input length increases, the performance first increases and then gradually decreases. The main reason is that a short sequence cannot provide sufficient temporal periodicity and trend information. But too large input length may introduce noises for future prediction, leading to performance degradation.\n\nThen, we vary $d$ from 8 to 128. The results are reported in Figure ~\\ref{hidden_units}. We can observe that the performance first increases and then remains stable. However, too large $d$ leads to extra computation overhead. Therefore, set the $d=64$ is enough to achieve a satisfactory result.\n\nAfter that, to test the impact of the distance threshold in HSG, we vary $\\epsilon$ from 6 to 20. The results are reported in Figure ~\\ref{epsilon}.\nAs $\\epsilon$ increases, we observe a performance first increase then slightly decrease, which is perhaps because a large $\\epsilon$ integrates too many stations, which introduces more noises during spatial autocorrelation modeling.\n\nFinally, we vary the prediction step $\\tau$ from 3 to 48. The results are reported in Figure ~\\ref{tau}. We observe the performance drops rapidly when $\\tau$ goes large. The reason perhaps is the uncertainty increases when $\\tau$ is large, and it is difficult for the machine learning model to quantify such uncertainty.\n\n\\subsection{Qualitative study}\n\\begin{figure}\n \\centering\n \n \\vspace{-4pt}\n \\includegraphics[width=8cm]{exp\/casestudy.pdf}\n \\vspace{-4pt}\n \\caption{Visualization of the learned attention weight of neighboring air quality and weather stations.}\\label{fig:casestudy}\n \\vspace{-7mm}\n\\end{figure}\n\nFinally, we qualitatively analyze why MasterGNN can yield better performance on both tasks. \nFigure~\\ref{fig:casestudy}~(a) shows the distribution of air quality and weather monitoring stations in the Haidian district, Beijing. \nWe perform a case study from 3:00 May 4, 2017, to 3:00 May 5, 2017.\nFirst, we take the air quality station $S_0$ as an illustrative example to show how weather stations help air quality prediction.\nIn MasterGNN, the attention score reflects the relative importance of each neighboring station for prediction.\nFigure~\\ref{fig:casestudy}~(b) visualizes the attention weights of eight neighboring weather stations of $S_0$.\nAs can be seen, the attention weights of weather stations $S_1$, $S_2$ and $S_4$ abruptly increase because a strong wind blows through this area, which results in notable air pollutant dispersion in the next five hours.\nClearly, the weather stations provide additional knowledge for air quality prediction.\nThen, we take the weather station $S_2$ to depict air quality stations' impact on weather prediction.\nFigure~\\ref{fig:casestudy}~(c) shows the attention weights of neighbouring station of $S_2$.\nWe observe the nearest air quality station $S_3$ is the most influential station during 5:00 May 4 2017 and 10:00 May 4 2017, while $S_7$ plays a more important role during 15:00 May 4 2017 and 21:00 May 4 2017, corresponding to notable air pollution concentration changes during these periods.\nThe above observations further validate the additional knowledge introduced by the joint modeling of heterogeneous stations for both predictive tasks.\n\n\\eat{\nFigure~\\ref{fig:casestudy} presents a local map of Beijing, several air quality and weather stations are scattered in geographical space. We perform a case study from 3:00 May 4, 2017 to 3:00 May 5, 2017. \nTo show how weather conditions impacts air quality, we take the air quality station $S_0$ as an example to visualize the attention weights with eight neighboring weather stations in Figure~\\ref{fig:casestudy}. The attention scores semantically represent the relative importance of each weather station. According to the attention weight matrix, the stations are not correlated before 22:00 May 4, 2017. Then, there came a strong wind and has lasted for 5 hours, the air quality dropped rapidly in this time period. We observe the attention weights increase when the wind blows strongly. Note that the nearest weather station $S_1$ have the largest attention scores, showing biggest impact to the selected air quality station. The attention weights gradually decreases with increasing distance. The above example have shown that the attention weights generated by our model are indeed meaningful in real world and can be easily interpreted.\nNext, we further illustrate how air quality impacts meteorology. As shown in Figure ~\\ref{fig:casestudy}, the attention weights of nearby air quality stations are higher than remote stations due to severe air pollution in this region. Another observation is that the attention weight of air quality to weather is lower than that of weather to air quality, this shows the influence of weather on air quality is much more significant.\nThe case study indicates that our model successfully captures the complex spatial correlation between the air quality stations and weather stations.\n}\n\\section{Introduction}\n\n\\begin{figure}\n \\centering\n \n \\vspace{-4pt}\n \\includegraphics[width=8cm]{fig\/distribution.pdf}\n \\vspace{-5pt}\n \\caption{\n Spatial distribution of air quality and weather monitoring stations in Beijing. Two types of stations are monitoring exclusively different but correlated air quality and weather conditions in different city locations.}\\label{dist}\n \\vspace{-7mm}\n\\end{figure}\n\n\n\nWith the rapid development of the economy and urbanization, people are increasingly concerned about emerging public health risks and environmental sustainability.\nAs reported by the World Health Organization~(WHO), air pollution is the world's largest environmental health risk~\\cite{campbell2019climate}, and the changing weather profoundly affects the local economic development and people's daily life~\\cite{baklanov2018urban}.\nThus, accurate and timely air quality and weather predictions are of great importance to urban governance and human livelihood.\nIn the past years, massive sensor stations have been deployed for monitoring air quality~(\\emph{e.g.},\\xspace PM2.5 and PM10) or weather conditions~(\\emph{e.g.},\\xspace temperature and humidity) and many efforts have been made for air quality and weather predictions~\\cite{liang2018geoman,yi2018deep,wang2019deep}.\n\nHowever, existing air quality and weather prediction methods are designated for either air quality or weather prediction, perhaps with one another as a side input~\\cite{zheng2015forecasting,yi2018deep}, but overlook the intrinsic connections and interactions between two tasks.\nDifferent from existing studies, this work is motivated by the following two insights. First, air quality and weather predictions are two highly correlated tasks and can be mutually enhanced.\nFor example, accurately predicting the regional wind condition can help model the future dispersion and transport of air pollutants~\\cite{ding2017air}.\nAs another example, modeling the future concentration of aerosol pollutants~(\\emph{e.g.},\\xspace PM2.5 and PM10) also can help predict the local climate~(\\emph{e.g.},\\xspace temperature and humidity), since aerosol elements can scatter or absorb the incoming radiation~\\cite{hong2020weakening}.\nSecond, as illustrated in Figure~\\ref{dist}, the geospatially distributed air quality and weather monitoring stations provide additional hints to improve both predictive tasks.\nThe air quality and weather condition variations in different city locations reflect the urban dynamics and can be exploited to improve spatiotemporal autocorrelation modeling.\n\nIn this work, we investigate the joint prediction of air quality and weather conditions by explicitly modeling the correlations and interactions between two predictive tasks.\nHowever, two major challenges arise in achieving this goal.\n(1) \\emph{Observation heterogeneity}.\nAs illustrated in Figure \\ref{dist}, the geo-distributed air quality and weather stations are heterogeneous spatial objects that are monitoring exclusively different atmospheric conditions. \nExisting methods~\\cite{zheng2013u,wang2019deep} are initially designed to model homogeneous spatial objects~(\\emph{i.e.},\\xspace either air quality or weather stations), which are not suitable for joint air quality and weather predictions.\nTherefore, the first challenge is how to capture spatiotemporal autocorrelations among heterogeneous monitoring stations to mutually benefit air quality and weather prediction.\n(2) \\emph{Compounding observation error vulnerability}.\nIn practice, observations reported by monitoring stations are often noisy due to the sensor error and environmental interference~\\cite{yi2016st}.\nHowever, most existing prediction models~\\cite{zheng2015forecasting,liang2018geoman,yi2018deep} rely on spatiotemporal dependency modeling between stations, which is susceptible to local perturbations and noise propagation~\\cite{zugner2018adversarial,bengio2015scheduled}.\nMore severely, jointly modeling spatiotemporal autocorrelation among air quality and weather monitoring stations will further accumulate errors from both spatial and temporal domains.\nAs a result, it is challenging to learn robust representations to resist compounding observation error for joint air quality and weather predictions.\n\nTo tackle the above challenges, we propose the \\emph{\\textbf{M}ulti-\\textbf{a}dversarial \\textbf{s}patio\\textbf{te}mporal \\textbf{r}ecurrent \\textbf{G}raph \\textbf{N}eural \\textbf{N}etworks}~(\\textbf{MasterGNN}) for robust air quality and weather joint predictions.\nSpecifically, we first devise a heterogeneous recurrent graph neural network to simultaneously incorporate spatial and temporal autocorrelations among heterogeneous stations conditioned on dynamic urban contextual factors~(\\emph{e.g.},\\xspace POI distributions, traffics).\nThen, we propose a multi-adversarial learning framework to against noise propagation from both microscopic and macroscopic perspectives.\nBy proactively simulating perturbations and maximizing the divergence between target and fake observations, multiple discriminators dynamically regularize air quality and weather station representations to resist the propagation of observation noises~\\cite{durugkar2016generative}.\nMoreover, we introduce a multi-task adaptive training strategy to improve the joint prediction performance by automatically balancing the importance of multiple discriminators.\nFinally, we conduct extensive experiments on two real-world datasets collected from Beijing and Shanghai, and the proposed model consistently outperforms seven baselines on both air quality and weather prediction tasks.\n\\eat{\nTo sum up, our main contributions are as follows:\n\\begin{itemize}\n\t\\item We study the novel air quality and weather co-prediction problem by exploiting the inner-connection between air quality and weather conditions.\n\t\\item We propose a heterogeneous recurrent graph neural network to capture the spatiotemporal autocorrelation among heterogeneous monitoring stations.\n\t\\item We introduce a multi-adversarial learning framework coupled with a multi-objective optimization strategy to resist observation noises.\n\t\\item We conduct extensive experiments on two real-world datasets collected from Beijing and Shanghai. The proposed model outperforms seven baselines by at least 10.7\\% and 5.5\\% on air quality and weather prediction tasks, respectively.\n\\end{itemize}\n}\n\n\\section{Methodology}\nMasterGNN simultaneously makes air quality and weather predictions based on the following intuitions.\n\n\\textbf{Intuition 1: Heterogeneous spatiotemporal autocorrelation modeling}. The geo-distributed air quality and weather monitoring stations provide additional spatiotemporal hints for air quality and weather prediction. The model should be able to simultaneously incorporate spatial and temporal autocorrelations of such heterogeneous monitoring stations for joint air quality and weather prediction.\n\n\\textbf{Intuition 2: Robust representation learning}. \nThe spatiotemporal autocorrelation based joint prediction model is vulnerable to observation noises and suffers from compounding propagation errors. The model should be robust to error propagation from both spatial and temporal domains.\n\n\\textbf{Intuition 3: Adaptive model training}. Learning robust representations introduces extra optimization objectives, which guiding the predictive model to approximate the underlying data distribution from different semantic aspects. The model should be able to balance the optimization of multiple objectives for optimal prediction adaptively. \n\n\\subsection{Framework Overview}\n\n\\begin{figure}[t]\n \\centering\n \n \n \\includegraphics[width=7.5cm]{fig\/framework.pdf}\n \\vspace{-0.1cm}\n \\caption{An overview of MasterGNN.}\\label{framework}\n \\vspace{-0.6cm}\n\\end{figure}\n\nFigure \\ref{framework} shows the architecture of our approach, including three major tasks: \n(i) modeling the spatiotemporal autocorrelation among air quality and weather stations for joint prediction; (ii) learning robust station representations via multi-adversarial learning; (iii) exploiting adaptive training strategy to ease the model learning.\nSpecifically, in the first task, we propose a \\emph{Heterogeneous Recurrent Graph Neural Network}~(HRGNN) to jointly incorporate spatial autocorrelation between heterogeneous monitoring stations and past-current temporal autocorrelation of each monitoring station.\nIn the second task, we develop a \\emph{Multi-Adversarial Learning} framework to resist the propagation of observation noises via (i) microscopic discriminators against adversarial attacks respectively from the spatial and temporal domain, and (ii) macroscopic discriminator against adversarial attacks from a global view of the city.\nIn the third task, we introduce a \\emph{Multi-Task Adaptive Training} strategy to automatically balance the optimization of multiple discriminative losses in multi-adversarial learning.\n\n\n\\subsection{Heterogeneous Recurrent Graph Neural Network}\n\n\n\\subsubsection{The base model.}\nWe adopt the encoder-decoder architecture~\\cite{zhang2020tkde,zhang2020semi} as the base model.\nSpecifically, the encoder step projects the historical observation sequence of all stations $\\mathbfcal{X}$ to a hidden state $\\mathbf{H}^t=f_{encoder}(W\\mathbfcal{X}+b)$, where $f_{encoder}(\\cdot)$ is parameterized by a neural network.\nIn the decoder step, we employ another neural network $f_{decoder}(\\cdot)$ to generate air quality and weather predictions $(\\hat{\\mathbf{Y}}^{t+1},\\hat{\\mathbf{Y}}^{t+2},\\dots,\\hat{\\mathbf{Y}}^{t+\\tau})=f_{decoder}(\\mathbf{H}^t, \\mathbf{C})$, where $\\tau$ is the future time steps we aim to predict, and $\\mathbf{C}$ are contextual features~\\cite{liu2020incorporating}.\nSimilar with conventional air quality or weather prediction models~\\cite{liang2018geoman,li2020weather}, the heterogeneous recurrent graph neural network aims to minimize the \\emph{Mean Square Error}~(MSE) between the target observations and predictions,\n\\begin{equation}\\label{equ:mse}\n\\mathcal{L}_g=\\frac{1}{\\tau}\\sum^{\\tau}_{i=1}||\\hat{\\mathbf{Y}}^{t+i}-\\mathbf{Y}^{t+i}||^2_2.\n\\end{equation}\n\n\\subsubsection{Incorporating spatial autocorrelation.}\nIn the spatial domain, the regional concentration of air pollutants and weather conditions are highly correlated and mutually influenced.\nInspired by the recent variants~\\cite{wang2019heterogeneous,zhang2019heterogeneous,liu2021vldb} of graph neural network on handling non-Euclidean semantic dependencies on heterogeneous graphs, we devise a context-aware heterogeneous graph attention block (CHAT) to model spatial interactions between heterogeneous stations.\nWe first construct the heterogeneous station graph to describe the spatial adjacency of each station.\n\n\\begin{myDef}\n\\textbf{Heterogeneous station graph}. A heterogeneous station graph (HSG) is defined as $\\mathcal{G}=\\{\\mathcal{V},\\mathcal{E}, \\psi \\}$, where $\\mathcal{V}=S$ is the set of monitoring stations, $\\psi$ is a mapping function indicates the type of each station, and $\\mathcal{E}$ is a set of directed edges indicating the connectivity among monitoring stations, defined as\n\\begin{equation}\ne_{ij}=\\left\\{\n\\begin{array}{lr}\n1,\\quad d_{ij}<\\epsilon\\\\\n0,\\quad otherwise\n\\end{array}\n\\right.,\n\\end{equation}\nwhere $d_{ij}$ is the spherical distance between station $s_i$ and station $s_j$, $\\epsilon$ is a distance threshold.\n\\end{myDef}\n\nDue to the heterogeneity of monitoring stations, there are two types of vertices~(air quality station $s^a$, weather station $s^w$) and four types of edges, \\emph{i.e.},\\xspace $\\Psi=\\{s^a$-$s^a$, $s^w$-$s^w$, $s^a$-$s^w$, $s^w$-$s^a\\}$. \nWe use $\\psi(i)$ to denote the station type of $s_i$, and $r \\in \\Psi$ to denote the semantic type of each edge.\n\nFormally, given an observation $\\mathbf{x}_i$ at a particular time step, we first devise a type-specific transformation layer to project heterogeneous observations into unified feature space,\n$\\widetilde{\\mathbf{x}}_i = \\mathbf{W}^{\\psi(i)} \\mathbf{x}_i$,\nwhere $\\widetilde{\\mathbf{x}}_i\\in\\mathcal{R}^d$ is a low-dimensional embedding vector, $\\mathbf{W}^{\\psi(i)}\\in\\mathcal{R}^{|\\mathbf{x}_i|\\times d}$ is a learnable weighted matrix shared by all monitoring stations of type $\\psi(i)$.\n\nThen, we introduce a type-dependent attention mechanism to quantify the non-linear correlation between homogeneous and heterogeneous stations under different contexts. Given a station pair $(s_i,s_j)$ which are connected by an edge of type $r$, the attention score is defined as\n\\begin{equation}\na^r_{ij} =\\frac{Attn(\\widetilde{\\mathbf{x}}_i,\\widetilde{\\mathbf{x}}_j,\\mathbf{c}_i,\\mathbf{c}_j,d_{ij})}{\\sum_{k\\in\\mathcal{N}^r_{i}}Attn(\\widetilde{\\mathbf{x}}_i,\\widetilde{\\mathbf{x}}_k,\\mathbf{c}_i,\\mathbf{c}_k,d_{ik})},\n\\end{equation}\nwhere $Attn(\\cdot)$ is a concatenation based attention function, $\\mathbf{c}_i,\\mathbf{c}_j\\in\\mathbf{C}$ are contextual features of station $s_i$ and $s_j$, $d_{ij}$ is the spherical distance between $s_i$ and $s_j$, and $\\mathcal{N}^r_i$ is the set of type-specific neighbor stations of $s_i$ in $\\mathcal{G}$.\nBased on $a^r_{ij}$, we define the context-aware heterogeneous graph convolution operation to update the type-wise station representation,\n\\begin{equation}\n\\widetilde{\\mathbf{x}}_i^{r\\prime} = GConv(\\widetilde{\\mathbf{x}}_{i}, r) = \\sigma(\\sum_{j\\in \\mathcal{N}^r_i}\\alpha^r_{ij}\\mathbf{W}^r\\widetilde{\\mathbf{x}}_j),\n\\end{equation}\nwhere $\\widetilde{\\mathbf{x}}^{r\\prime}_{i}$ is the aggregated node representation based on edge type $r$, $\\sigma$ is a non-linear activation function, $\\mathbf{W}^r\\in\\mathcal{R}^{d \\times d}$ is a learnable weighted matrix shared over all edges of type $r$.\nFinally, we obtain the updated station representation of $s_i$ by concatenating type-specific representations,\n\\begin{equation}\n\\widetilde{\\mathbf{x}}^{\\prime}_i = ||_{r\\in\\Psi}GConv(\\widetilde{\\mathbf{x}}_{i}, r),\n\\end{equation}\nwhere $||$ is the concatenation operation.\nNote that we can stack $l$ graph convolution layers to capture the spatial autocorrelation between $l$-hop heterogeneous stations.\n\n\\eat{\nConsider an air quality observation $\\mathbf{x}^a_i$ and a weather observation $\\mathbf{x}^w_j$ at a particular time step, we first employ a type-specific linear transformation layer to project heterogeneous observations into a unified feature space,\n\\begin{equation}\n\\mathbf{h}^a_i = \\mathbf{W}^a \\mathbf{x}^a_i,~~\\mathbf{h}^w_j = \\mathbf{W}^w \\mathbf{x}^w_j,\n\\end{equation}\nwhere $\\mathbf{h}^a_i\\in\\mathcal{R}^d$ and $\\mathbf{h}^w_j\\in\\mathcal{R}^d$ are low-dimensional embedding vectors, $\\mathbf{W}^a\\in\\mathcal{R}^{|\\mathbf{x}^a_i|\\times d}$ and $\\mathbf{W}^w\\in\\mathcal{R}^{|\\mathbf{x}^w_j|\\times d}$ are learnable weighted matrices shared by air quality and weather monitoring stations, respectively.\n\nWe use $\\psi_i$ to denote the node type of $v_i$, and $\\psi_{ij}$ to denote the edge type from $v_i$ to $v_j$.\n\nConsider an atmospheric observation $\\mathbf{x}_i$ at a particular time step, we first employ a type-specific linear transformation layer to project heterogeneous observations into a unified feature space,\n\\begin{equation}\n\\mathbf{h}_i = \\mathbf{W}_{\\psi_i} \\mathbf{x}_i,\n\\end{equation}\nwhere $\\mathbf{h}_i\\in\\mathcal{R}^d$ is a low-dimensional embedding vector, $\\mathbf{W}_{\\psi_i}\\in\\mathcal{R}^{|\\mathbf{X}_i|\\times d}$ is a learnable weighted matrix shared by all same-typed monitoring stations.\n}\n\n\\subsubsection{Incorporating temporal autocorrelation.}\nIn the temporal domain, the air quality and weather conditions also depend on previous observations.\nWe further extend the Gated Recurrent Units (GRU), a simple variant of recurrent neural network~(RNN), to integrate the temporal autocorrelation among heterogeneous observation sequences.\nConsider a station $s_i$, given the learned spatial representation $\\widetilde{\\mathbf{x}}^t_i$ at time step $t$, we denote the hidden state of $s_i$ at $t-1$ and $t$ as $\\mathbf{h}^{t-1}_i$ and $\\mathbf{h}^t_i$, respectively. The temporal autocorrelation between $\\mathbf{h}^{t-1}_i$ and $\\mathbf{h}^t_i$ is modeled by\n\\begin{equation}\n\\left\\{\n\\begin{aligned}\n &\\mathbf{h}^t_i=(1-\\mathbf{z}^t_i)\\odot \\mathbf{h}^{t-1}_i+\\mathbf{z}^t_i\\odot \\widetilde{\\mathbf{h}}^t_i\\\\\n &\\mathbf{r}^{t}_i = \\sigma{(\\mathbf{W}^{\\psi(i)}_r[\\mathbf{h}^{t-1}_i\\parallel\\widetilde{\\mathbf{x}}^{t}_i]+\\mathbf{b}^{\\psi(i)}_r)}\\\\\n &\\mathbf{z}^{t}_i = \\sigma(\\mathbf{W}^{\\psi(i)}_z[\\mathbf{h}^{t-1}_i\\parallel\\widetilde{\\mathbf{x}}^{t}_i]+\\mathbf{b}^{\\psi(i)}_z)\\\\\n &\\widetilde{\\mathbf{h}}^{t}_i = \\tanh(\\mathbf{W}^{\\psi(i)}_{\\widetilde{h}}[\\mathbf{r}^{t}_i\\odot\\mathbf{h}^{t-1}_i\\parallel\\widetilde{\\mathbf{x}}^t_i]+\\mathbf{b}^{\\psi(i)}_{\\widetilde{h}})\n\\end{aligned},\n\\right.\n\\end{equation}\nwhere $\\mathbf{o}^t_i$, $\\mathbf{z}^t_i$ denote reset gate and update gate at time stamp t, $W_{o}^{\\psi(i)}$, $W_{z}^{\\psi(i)}$, $W_{\\widetilde{h}}^{\\psi(i)}$, $b_{o}^{\\psi(i)}$, $b_{z}^{\\psi(i)}$, $b_{\\widetilde{h}}^{\\psi(i)}$ are trainable parameters shared by all same-typed monitoring stations, and $\\odot$ represents the Hadardmard product. \n\n\\subsection{Multi-Adversarial Learning}\n\n\\begin{figure}[t]\n \\centering\n \n \n \\includegraphics[width=8.5cm]{fig\/adver_train.pdf}\n \\vspace{-0.5cm}\n \\caption{Discriminators in multi-adversarial learning.}\\label{adver_train}\n \\vspace{-0.7cm}\n\\end{figure}\n\nWe further introduce the multi-adversarial learning framework to obtain robust station representations.\nBy proactively simulating and resisting noisy observations via a \\emph{minmax} game between the generator and discriminator~\\cite{goodfellow2014generative}, adversarial learning can make the generator more robust and generalize better to the underlying true distribution $p_{true}(\\mathbf{Y} | \\mathbfcal{X}, \\mathbf{C})$.\nCompared with standard adversarial learning, our multi-adversarial learning framework encourages the generator to make more consistent predictions from both spatial and temporal domains.\nSpecifically, we function the HRGNN $\\mathbf{\\phi}$ as the generator $\\hat{\\mathbf{Y}}=G(\\mathbfcal{X}, \\mathbf{C};\\mathbf{\\phi})$ for joint predictions.\nMoreover, we propose a set of distinct discriminators $\\{D_k(\\cdot;\\theta_k)\\}^K_{k=1}$ to distinguish the real observations and predictions from both microscopic and macroscopic perspectives, as detailed below.\n\n\\subsubsection{Microscopic discriminators.}\nMicroscopic discriminators aims at enforcing the generator to approximate the underlying spatial and temporal distributions, \\emph{i.e.},\\xspace pair-wise spatial autocorrelation and stepwise temporal autocorrelation.\n\n\\emph{Spatial discriminator}. From the spatial domain, the adversarial perturbation induces the compounding propagation error through the HSG, as illustrated in Figure~\\ref{adver_train}~(a). Consider a time step $t$, the spatial discriminator $D_s(\\mathbf{y}^t;\\theta_s)$ aims to maximize the accuracy of distinguishing city-wide real observations and predictions at the current time step,\n\\begin{equation}\\label{equ:dis1}\n\\mathcal{L}_{s}=\\log D_s(\\mathbf{y}^t;\\theta_s)+log(1-D_s(\\hat{\\mathbf{y}}^t;\\theta_s)),\n\\end{equation}\nwhere $D_s(\\mathbf{y}^t;\\theta_s)$ is parameterized by the context-aware heterogeneous graph attention block in HRGNN followed by a multi-layer perception, $\\mathbf{y}^t$ and $\\hat{\\mathbf{y}}^t$ are ground truth observations and predicted conditions of the city, respectively.\n\n\\emph{Temporal discriminator.} From the temporal domain, the adversarial perturbation induces accumulated error for each station from previous time steps. \nAs depicted in Figure~\\ref{adver_train}~(b), the temporal discriminator $D_t(\\mathbf{y}_i;\\theta_t)$ outputs a probability indicating how likely a observation sequence $\\mathbf{y}_i$ of a particular station $s_i$ is from the real data distribution $p_{true}(\\mathbf{y}_i | \\mathbf{x}_i, \\mathbf{c}_i)$ rather than the generator $p_{\\phi}(\\mathbf{y}_i | \\mathbf{x}_i, \\mathbf{c}_i)$.\n\\begin{equation}\\label{equ:dis2}\n\\mathcal{L}_{t}=\\log D_t(\\mathbf{y}_i;\\theta_t)+log(1-D_t(\\hat{\\mathbf{y}}_i;\\theta_t)).\n\\end{equation}\nDifferent from the spatial discriminator, $D_t(\\mathbf{y}_i;\\theta_t)$ is parameterized by the temporal block in HRGNN followed by a multi-layer perception. \n$\\mathbf{y}_i$ is a sequence of observations in past $T$ and future $\\tau$ time steps.\n\n\\subsubsection{Macroscopic discriminator.}\nAs illustrated in Figure~\\ref{adver_train}~(c), we further propose a macroscopic discriminator to capture the globally underlying distribution, denoted by $D_m(\\mathbf{Y};\\theta_m)$. In particular, the macroscopic discriminator aims to maximize the accuracy of distinguishing the ground truth observations and predictions generated from $G$ from a global view,\n\\begin{equation}\\label{equ:dis3}\n\\mathcal{L}_{m}=\\log D_m(\\mathbf{Y};\\theta_m)+log(1-D_m(G(\\mathbfcal{X},\\mathbf{C};\\mathbf{\\phi});\\theta_m)),\n\\end{equation}\nwhere $D_m(\\mathbf{Y};\\theta_m)$ is parameterized by a GRU followed by a multi-layer perception.\nNote that for efficiency concern, the input of $D_m(\\mathbf{Y};\\theta_m)$ is a simple concatenation of observations from all stations $S$ in past $T$ and future $\\tau$ time steps, but other graph convolution operations are also applicable.\n\n\\eat{\nGiven the generator $G$ and discriminators $\\{D_i\\}^K_{i=1}$,\nthe minmax function of the multi-adversarial learning is\n\\begin{equation}\n\\begin{aligned}\n\\min_{G} \\max_{\\{D_i\\}^K_{i=1}} V(\\{D_i\\}^K_{i=1}&,G)=\\mathbb{E}_{\\mathbf{Y} \\sim p_{\\emph{true}}}[\\log D_i(\\mathbf{Y};\\theta_i)]+\\\\\n&\\mathbb{E}_{\\mathbfcal{X} \\sim p_{\\phi}}[\\log (1-D_i(G(\\mathbfcal{X},\\mathbf{C};\\phi)))].\n\\end{aligned}\n\\end{equation}\n}\n\n\\subsection{Multi-Task Adaptive Training}\nIt is widely recognized that adversarial training suffers from the unstable and mode collapse problem, where the generator is easy to over-optimized for a particular discriminator~\\cite{salimans2016improved,neyshabur2017stabilizing}.\nTo stabilize the multi-adversarial learning and achieve overall accuracy improvement, we introduce a multi-task adaptive training strategy to dynamically re-weight discriminative loss and enforce the generator to perform well in various spatiotemporal aspects.\n\nSpecifically, the objective of MasterGNN is to minimize \n\\begin{equation}\\label{equ:all-loss}\n\\mathcal{L}=\\mathcal{L}_g+\\sum^K_{i=1}\\lambda_i \\mathcal{L}_{d_i},\n\\end{equation}\nwhere $\\mathcal{L}_g$~(Equation~\\ref{equ:mse}) is the predictive loss, $\\mathcal{L}_D=\\{\\mathcal{L}_{d_i}\\}^K_{i=1}$ (Equation~\\ref{equ:dis1}-\\ref{equ:dis3}) are discriminative losses, and\n$\\lambda_i$ is the importance of the corresponding discriminative loss.\n\nSuppose $\\mathbf{H}_{i}$ and $\\hat{\\mathbf{H}}_{i}$ are intermediate hidden states in discriminator $D_i$ based on real observations $\\mathbf{Y}$ and predictions $\\hat{\\mathbf{Y}}$.\nWe measure the divergence between $\\mathbf{H}_{i}$ and $\\hat{\\mathbf{H}}_{i}$ by $\\gamma_i=sim(\\sigma(\\mathbf{H}_{i}),\\sigma(\\hat{\\mathbf{H}}_{i}))$,\nwhere $sim(\\cdot,\\cdot)$ is an Euclidean distance based similarity function, $\\sigma$ is the sigmoid function. \nIntuitively, $\\gamma_i$ reflects the hardness of $D_i$ to distinguish the real sample.\nIn each iteration, we re-weight discriminative losses by \n$\\lambda_i = \\frac{exp(\\gamma_i)}{\\sum^{K}_{k=1}exp(\\gamma_k)}$.\nIn this way, the generator pays more attention to discriminators with a larger room to improve, and will result in better prediction performance.\n\n\n\\eat{\nSpecifically, for the generative loss $\\mathcal{L}_g$ and $K$ discriminative losses $\\mathcal{L}_D=[\\mathcal{L}_1,\\mathcal{L}_2,...,\\mathcal{L}_K]$,\nwhere each $\\mathcal{L}_k=\\mathbb{E}_{\\mathbfcal{X} \\sim p_{\\phi}}[\\log(1-D_i(G(\\mathbfcal{X},\\mathbf{C};\\phi)))]$ is the loss provided by the i-th discriminator.\nWe use a dynamic weight $\\lambda_i$ to control the importance of discriminator $i$. A multi-task training strategy is designed to adaptively adjust $\\lambda_i$. Suppose $h_{i,f}$ and $h_{i,r}$ are the hidden state of discriminator $i$ derived from future predictions and ground truth observations, respectively. We first use a Sigmoid function to map the hidden state to the range [0,1]. Then we leverage a function to measure the similarity between $h_{i,f}$ and $h_{i,r}$, defined as follows\n\\begin{equation}\n\\gamma_i=sim(h_{i,f},h_{i,r}).\n\\end{equation}\nIn this paper, we use Euclidean Distance as similarity function. When $\\gamma_i$ becomes large, it is intuitive to conclude the discriminator $i$ can easily distinguish the predictions and ground truth observations, which means the generator performs poorly on discriminator $i$. We directly adopt $\\gamma_i$ as the weight to focus on gradients from discriminator $i$.\nWe employ a softmax function to convert the dynamic weights into a probability distribution.\n\\begin{equation}\n\\lambda_i = \\frac{exp(\\gamma_i)}{\\sum^{K}_{k=1}exp(\\gamma_k)}.\n\\end{equation}\nDuring training, our proposed strategy dynamically adjusts the influence of various discriminators. \nUpdates of generator are performed to minimize the following loss\n\\begin{equation}\n\\mathcal{L}=\\mathcal{L}_g+\\sum^K_{i=1}\\lambda_i \\mathcal{L}_{i},\n\\end{equation}\n\nWe update the multiple discriminators while keeping the parameters of the generator fixed. Each discriminator is minimized using the loss described in Eq.12.\n}\n\n\n\n\n\n\\section{Preliminaries}\n\n\n\nConsider a set of monitoring stations $S=S^a\\cup S^w$, where $S^a=\\{s^a_i\\}^m_{i=1}$ and $S^w=\\{s^w_j\\}^n_{j=1}$ are respectively air quality and weather station sets.\nEach station $s_i \\in S$ is associated with a geographical location $l_i$~(\\emph{i.e.},\\xspace latitude and longitude) and a set of time-invariant contextual features $\\mathbf{c}_i \\in \\mathbf{C}$.\n\n\\begin{myDef}\n\\textbf{Observations}. Given a monitoring station $s_i \\in S$, the observations of $s_i$ at time step $t$ are defined as $\\mathbf{x}^{t}_i$, which is a vector of air-quality or weather conditions, depending on the station type. \n\\end{myDef}\n\nNote that the observations of two types of monitoring stations are different~(\\emph{e.g.},\\xspace PM2.5 and CO in air quality stations, while temperature and humidity in weather stations) and the observation dimensionality of the air quality station and the weather station are also different. \nWe use $\\mathbfcal{X}=\\{\\mathbf{X}^{a,t}\\}^T_{t=1} \\cup \\{\\mathbf{X}^{w,t}\\}^T_{t=1}$ to denote time-dependent observations of all stations in a time period $T$, use $\\mathbf{X}^{a,t}=\\{\\mathbf{x}^{a,t}_1, \\mathbf{x}^{a,t}_2, \\dots, \\mathbf{x}^{a,t}_{m}\\}$ and $\\mathbf{X}^{w,t}=\\{\\mathbf{x}^{w,t}_1, \\mathbf{x}^{w,t}_2, \\dots, \\mathbf{x}^{w,t}_{n}\\}$ to respectively denote observations of air quality and weather stations at time step $t$. We use $\\mathbf{X}_i=\\{\\mathbf{x}^{1}_i, \\mathbf{x}^{2}_i, \\dots, \\mathbf{x}^{T}_i\\}$ to denote all observations of station $s_i$.\nIn the following, without ambiguity, we will omit the superscript and subscript.\n\n\\eat{\n\\begin{myDef}\n\\textbf{Heterogeneous Station Graph}. A heterogeneous station graph (HSG) is defined as $\\mathcal{G}=\\{\\mathcal{V},\\mathcal{E}, \\phi \\}$, where $\\mathcal{V}=S$ is the set of monitoring stations, $\\phi$ is a mapping set indicates the type of each station, and $\\mathcal{E}$ is a set of edges indicating the connectivity among monitoring stations, defined as\n\\begin{equation}\ne_{ij}=\\left\\{\n\\begin{array}{lr}\n1,\\quad dist(s_i,s_j)<\\epsilon\\\\\n0,\\quad otherwise\n\\end{array}\n\\right.,\n\\end{equation}\nwhere $dist(\\cdot,\\cdot)$ is the spherical distance between station $s_i$ and station $s_j$, $\\epsilon$ is a distance threshold.\n\\end{myDef}\n}\n\n\\begin{problem}\n\\textbf{Joint air quality and weather predictions}. \nGiven a set of monitoring stations $S$, contextual features $\\mathbf{C}$, and historical observations $\\mathbfcal{X}$, our goal at a time step $t$ is to simultaneously predict air quality and weather conditions for all $s_i \\in S$ over the next $\\tau$ time steps,\n\\begin{equation}\n(\\hat{\\mathbf{Y}}^{t+1},\\hat{\\mathbf{Y}}^{t+2},...,\\hat{\\mathbf{Y}}^{t+\\tau}) \\leftarrow \\mathcal{F}(\\mathbfcal{X}, \\mathbf{C}),\n\\end{equation}\nwhere $\\hat{\\mathbf{Y}}^{t+1}=\\hat{\\mathbf{Y}}^{a,t+1} \\cup \\hat{\\mathbf{Y}}^{w,t+1}$ is the estimated target observations of all stations at time step $t+1$, and $\\mathcal{F}(\\cdot)$ is the mapping function we aim to learn. Note that $\\hat{\\mathbf{y}}^{t+1}_i \\in \\hat{\\mathbf{Y}}^{t+1}$ is also station type dependent, \\emph{i.e.},\\xspace air quality or weather observations for corresponding stations.\n\\end{problem}\n\n\\section{Related work}\n\n\\textbf{Air quality and weather prediction}.\nExisting literature on air quality and weather prediction can be categorized into two classes.\n(1) \\emph{Numerical-based models} make predictions by simulating the dispersion of various air quality or weather elements based on physical laws~\\cite{lorenc1986analysis,vardoulakis2003modelling,liu2005prediction,richardson2007weather}.\n(2) \\emph{Learning-based models} utilize end-to-end machine learning methods to capture spatiotemporal correlations based on historical observations and various urban contextual data~(\\emph{e.g.},\\xspace POI distributions, traffics)~\\cite{chen2011comparison,zheng2013u,cheng2018neural}.\nRecently, many deep learning models~\\cite{yi2018deep,wang2019deep,lin2018exploiting} have been proposed to enhance the performance of air quality and weather prediction.\nBy leveraging the representation capacity of deep learning for spatiotemporal autocorrelation modeling, learning-based models usually achieve better prediction performance than numerical-based models.\nUnlike the above approaches, our method explicitly models the correlations and interactions between the air quality and weather prediction task and achieves a better performance.\n\n\n\\textbf{Adversarial Learning}.\nAdversarial learning~\\cite{goodfellow2014generative} is an emerging learning paradigm for better capturing the data distribution via a minmax game between a generator and a discriminator.\nIn the past years, adversarial learning has been widely applied to many real-world application domains, such as sequential recommendation~\\cite{yu2017seqgan} and graph learning~\\cite{Wang2018GraphGANGR}.\nRecently, we notice several multi-adversarial frameworks have been proposed for improving image generation tasks~\\cite{nguyen2017dual,hoang2018mgan,albuquerque2019multi}.\nInspired by the above studies, we extend the multi-adversarial learning paradigm to the environmental science domain and introduce an adaptive training strategy to improve the stability of adversarial learning.\n\n\\eat{\nleverages a deep distributed fusion network for air quality prediction. The study in ~\\cite{grover2015deep} model the statistics of a set of weather-related variables via a hybrid approach based on deep neural networks. More recently, DUQ~\\cite{wang2019deep} introduced a deep uncertainty quantification method for weather forecasting.\n\nare widely used in the earlier study, such as ARIMA \\cite{chen2011comparison}, SVM~\\cite{wang2014research}, and artificial neural networks~\\cite{brunelli2007two}, they fail to capture spatiotemporal dynamics from the original data. With the development of deep learning techniques, many deep learning models are proposed to solve air quality and weather forecasting tasks. For example, DeepAir~\\cite{yi2018deep} leverages a deep distributed fusion network for air quality prediction. The study in ~\\cite{grover2015deep} model the statistics of a set of weather-related variables via a hybrid approach based on deep neural networks. More recently, DUQ~\\cite{wang2019deep} introduced a deep uncertainty quantification method for weather forecasting.\n\n\n\\subsection{Graph neural network.}\nRecent years deep learning on graphs has gained much attention, which promotes the development of graph neural networks~\\cite{kipf2016semi}. For each node on the graph, graph neural networks learn a function to aggregate features of its neighbours and generate new node embedding, which encodes both feature information and local graph structure. For example, graph attention network (GAT)~\\cite{velivckovic2017graph} uses self-attention mechanism to select important neighbours adaptively and then aggregate them with different weights. Besides, researchers also extend graph neural networks on heterogeneous graphs~\\cite{zhang2019heterogeneous}. Graph neural networks are widely used in many applications, such as traffic flow prediction~\\cite{li2017diffusion}, ride-hailing demand forecasting~\\cite{geng2019spatiotemporal} and parking availability prediction~\\cite{zhang2020semi}.\n}\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} 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\\\\}\n\\newcommand{\\varepsilon}{\\varepsilon}\n\n\\numberwithin{equation}{section}\n\\newcommand{{\\nabla}}{{\\nabla}}\n\\newcommand{{\\mathrm{c.c.}}}{{\\mathrm{c.c.}}}\n\n\\newcommand{{\\veps}}{{\\varepsilon}}\n\n\n\\newcommand{\\mathsf{Sp}}{\\mathsf{Sp}}\n\\newcommand{\\mathsf{SU}}{\\mathsf{SU}}\n\\newcommand{\\mathsf{SL}}{\\mathsf{SL}}\n\\newcommand{\\mathsf{GL}}{\\mathsf{GL}}\n\\newcommand{\\mathsf{SO}}{\\mathsf{SO}}\n\\newcommand{\\mathsf{U}}{\\mathsf{U}}\n\\newcommand{\\mathsf{S}}{\\mathsf{S}}\n\\newcommand{\\mathsf{PSU}}{\\mathsf{PSU}}\n\\newcommand{\\mathsf{PSL}}{\\mathsf{PSL}}\n\\newcommand{\\mathsf{OSp}}{\\mathsf{OSp}}\n\\newcommand{\\mathsf{Spin}}{\\mathsf{Spin}}\n\\newcommand{\\mathsf{Mat}}{\\mathsf{Mat}}\n\\newcommand{\\mathsf{ISO}}{\\mathsf{ISO}}\n\n\n\\begin{document}\n\n\\begin{titlepage}\n\\begin{flushright}\nOctober, 2020 \\\\\n\\end{flushright}\n\\vspace{5mm}\n\n\\begin{center}\n{\\Large \\bf \nMassless particles in five and higher dimensions}\n\\end{center}\n\n\\begin{center}\n\n{\\bf Sergei M. Kuzenko and Alec E. Pindur} \\\\\n\\vspace{5mm}\n\n\\footnotesize{\n{\\it Department of Physics M013, The University of Western Australia\\\\\n35 Stirling Highway, Perth W.A. 6009, Australia}}\n~\\\\\n\\vspace{2mm}\nEmail: \\texttt{ \nsergei.kuzenko@uwa.edu.au, 21504287@student.uwa.edu.au}\\\\\n\n\n\\end{center}\n\n\\begin{abstract}\n\\baselineskip=14pt\nWe describe a five-dimensional analogue of Wigner's operator equation \n${\\mathbb W}_a = \\lambda P_a$, where ${\\mathbb W}_a $ is the Pauli-Lubanski vector, \n$P_a$ the energy-momentum operator, and $\\lambda$ the helicity of a massless particle. \nHigher dimensional generalisations are also given.\n\\end{abstract}\n\\vspace{5mm}\n\n\n\n\\vfill\n\n\\vfill\n\\end{titlepage}\n\n\\newpage\n\\renewcommand{\\thefootnote}{\\arabic{footnote}}\n\\setcounter{footnote}{0}\n\n\n\\allowdisplaybreaks\n\n\n\\section{Introduction}\n\nThe unitary representations of the Poincar\\'e group in four dimensions were classified \nby Wigner in 1939 \\cite{Wigner}, see \\cite{Weinberg95} for a recent review. Our modern understanding of elementary particles is based on this classification. \n\nUnitary representations of the Poincar\\'e group $\\mathsf{ISO}_0(d-1, 1)$ in higher dimensions,\n$d>4$, have been studied in the literature, see, e.g., \\cite{BB}. However, there still remain some aspects that are not fully understood, see, e.g., \\cite{Weinberg2020}\nfor a recent discussion. In this note we analyse the \nirreducible massless representations of $\\mathsf{ISO}_0(4,1)$ with a finite (discrete) spin.\n\nWe recall that the Poincar\\'e algebra $\\mathfrak{iso} (d-1,1) $ in $d$ dimensions is characterised by the commutation relations\\footnote{We make use of the mostly plus Minkowski metric $\\eta_{ab}$ and normalise the Levi-Civita tensor $\\ve_{a_1 \\dots a_d}$ \nby $\\ve_{0 1\\dots d-1}=1$.} \n\\begin{subequations} \\label{PA}\n\\begin{eqnarray}\n \\big[P_{a}, P_{b}\\big] & = & 0 ~, \\\\\n \\big[J_{ab}, P_{c}\\big] & = & {\\rm i}\\eta_{ac}P_{b} - {\\rm i}\\eta_{bc}P_a ~,\\\\\n \\big[J_{a b}, J_{c d}\\big] & = & {\\rm i} \\eta_{a c}J_{b d} - {\\rm i} \\eta_{a d} J_{b c} + {\\rm i} \\eta_{b d} J_{a c} - {\\rm i} \\eta_{bc} J_{a d} ~.\n\\end{eqnarray}\n\\end{subequations}\nIn any unitary representation of (the universal covering group of) the Poincar\\'e group, \nthe energy-momentum operator\n$P_a$ and the Lorentz generators $J_{ab}$ are Hermitian.\nFor every dimension $d$, the operator $P^a P_a$ is a Casimir operator. Other Casimir operators are dimension dependent. \n\nIn four dimensions, the second Casimir operator is \n${\\mathbb W}^a{\\mathbb W}_a $,\nwhere \n\\begin{equation}\n{\\mathbb W}^a=\\frac12 \\ve^{abcd} { J}_{bc}{ P}_d\n\\label{PLv}\n\\end{equation}\nis the Pauli-Lubanski vector.\nUsing the commutation relations \n \\eqref{PA}, it follows that \n the Pauli-Lubanski vector is translationally invariant, \n\\begin{subequations}\n\\begin{eqnarray}\n\\big[ P_a, {\\mathbb W}_b \\big]=0~,\n\\end{eqnarray}\nand possesses the following properties:\n\\begin{eqnarray}\n{\\mathbb W}^a{ P}_a&=&0 ~, \\\\\n\\big[{ J}_{ab},{\\mathbb W}_c \\big]&=&{\\rm i} \\eta_{ac}{\\mathbb W}_b\n-{\\rm i} \\eta_{bc}{\\mathbb W}_a~,\n\\\\ \n\\big[ {\\mathbb W}_a,{\\mathbb W}_b \\big]&=& {\\rm i} \\ve_{abcd}{\\mathbb W}^c{P}^d~. \n\\end{eqnarray}\n\\end{subequations}\nThe irreducible massive representations are characterised by the conditions \n\\begin{subequations}\n\\begin{eqnarray}\n { P}^a{ P}_a &=&-m^2\\,{\\mathbbm 1}~, \\qquad m^2 > 0~, \n \\qquad \n {\\rm sign } \\,{ P}^0 >0~, \\label{1.4a} \\\\\n {\\mathbb W}^a{\\mathbb W}_a &=& m^2s(s+1)\\,{\\mathbbm 1}~, \n \\end{eqnarray} \n \\end{subequations}\nwhere the quantum number $s$ is called spin. Its possible values \nin different representations are \n$s= 0, 1\/2, 1, 3\/2 \\dots$. \nThe massless representations are characterised by the condition $P^aP_a =0$. \nFor the physically interesting massless representations, it holds that \n\\begin{eqnarray}\n{\\mathbb W}_a = \\lambda P_a~,\n\\label{1.5}\n\\end{eqnarray}\nwhere the parameter $\\lambda$ determines the representation and is called the helicity. Its possible values are $0, \\pm\\frac12, \\pm 1, $ and so on. The parameter $|\\lambda |$ is called the spin of a massless particle.\n\nIn this paper we present a generalisation of Wigner's equation \\eqref{1.5}\nto five and higher dimensions.\n\n\n\\section{Unitary representations of $\\mathsf{ISO}_0(4,1)$}\n\nThe five-dimensional analogue of \\eqref{PLv} is the Pauli-Lubanski tensor\n\\begin{equation}\n\\mathbb{W}^{ab} = \\frac{1}{2}\\ve^{abcde}J_{cd}P_e~.\n\\label{2.1}\n\\end{equation}\nIt is translationally invariant, \n\\begin{eqnarray}\n\\big[ {\\mathbb W}_{ab} , P_c \\big] =0~,\n\\end{eqnarray}\nand possesses the following properties:\n\\begin{subequations}\n\\begin{eqnarray}\n\\mathbb{W}_{ab} P^b &=& 0~, \\\\\n\\big[\\mathbb{W}_{ab}, J_{cd}\\big] &=& {\\rm i} \\eta_{ac}\\mathbb{W}_{bd} \n- {\\rm i} \\eta_{ad}\\mathbb{W}_{bc} - {\\rm i} \\eta_{bc}\\mathbb{W}_{ad} + {\\rm i} \\eta_{bd}\\mathbb{W}_{ac}~, \\\\\n \\big[\\mathbb{W}_{ab}, \\mathbb{W}_{cd}\\big] \n &=& {\\rm i}\\ve_{acdfg} {\\mathbb{W}_b}^f P^{g} - {\\rm i} \\ve_{bcdfg} {\\mathbb{W}_a}^f P^{g}~.\n\\end{eqnarray}\n\\end{subequations}\nMaking use of ${\\mathbb W}_{ab}$ allows one to construct two Casimir operators, which are \n\\begin{equation}\n\\mathbb{W}_{ab} \\mathbb{W}^{ab}~, \\qquad \n\\mathbb{H} := \\mathbb{W}^{ab}J_{ab}~.\n\\label{2.4}\n\\end{equation}\n\n\n\\subsection{Irreducible massive representations} \n\nThe irreducible massive representations of the Poincar\\'e group $\\mathsf{ISO}_0(4,1)$\nare characterised by two conditions\n\\begin{subequations}\\label{2.5}\n\\begin{eqnarray}\n \\frac{1}{8} \\Big( \\mathbb{W}^{ab}\\mathbb{W}_{ab} + m \\mathbb{H} \\Big) &=& \n m^2 s_1(s_1 + 1)\\mathbbm{1} ~, \\\\\n\\frac{1}{8} \\Big( \\mathbb{W}^{ab}\\mathbb{W}_{ab} - m\\mathbb{H} \\Big) &=& \nm^2 s_2(s_2 + 1)\\mathbbm{1} ~,\n\\end{eqnarray}\n\\end{subequations}\nin addition to \\eqref{1.4a}.\nHere $s_1$ and $s_2$ are two spin values corresponding to the two $\\mathsf{SU}(2)$ subgroups of the universal covering group \n$\\mathsf{Spin}(4) \\cong \\mathsf{SU}(2) \\times \\mathsf{SU}(2)$\nof the little group.\\footnote{The equations \\eqref{2.5} were independently derived during the academic year 1992-93 by Arkady Segal and David Zinger, who were undergraduates at Tomsk State University at the time.}\n\n\n\\subsection{Irreducible massless representations} \\label{section2.2}\n\nIt turns out that all irreducible massless representations of \n$\\mathsf{ISO}_0(4,1)$ with a finite spin are characterised by the condition \n\\begin{eqnarray}\n\\varepsilon_{abcde}P^{c}\\mathbb{W}^{de} =0 \\quad \\Longleftrightarrow \\quad\nP^{[a} {\\mathbb W}^{bc]} =0~.\n\\label{main}\n\\end{eqnarray}\nBoth Casimir operators \\eqref{2.4} are equal to zero in these representations, \n$\\mathbb{W}_{ab} \\mathbb{W}^{ab}=0$ and \n$ \\mathbb{W}^{ab}J_{ab} =0$.\n\n\nLet $ \\ket{p, \\sigma} $ be an orthonormal basis in the Hilbert space of one-particle states, where $p^a $ \ndenotes the momentum of a particle, $ P^a \\ket{p, \\sigma} = p^a \\ket{p, \\sigma} $,\nand $\\sigma$ stands for the spin degrees of freedom.\nFor a massless particle, we choose as our standard 5-momentum $k^{a} = (E,0,0,0,E)$. On this eigenstate:\n\\begin{equation}\n\\mathbb{W}^{ab} \\ket{k, \\sigma} = \\frac{1}{2}\\varepsilon^{abcde}J_{cd}P_{e} \\ket{k, \n\\sigma} = \\frac{E}{2} \\left(\\varepsilon^{abcd4}J_{cd} - \\varepsilon^{abcd0}J_{cd} \\right) \\ket{k, \\sigma}~.\n\\end{equation}\nRunning through the elements of $\\mathbb{W}^{ab}$, one finds:\n\\begin{equation} \\label{eq:eq5}\n\\begin{split}\n\\mathbb{W}^{01} = \\mathbb{W}^{41} = -EJ_{23} ~,\\qquad & \\mathbb{W}^{12} = E(J_{30} +J_{34}) ~,\\\\\n\\mathbb{W}^{02} = \\mathbb{W}^{42} = -EJ_{31} ~,\\qquad & \\mathbb{W}^{23} = E(J_{10} +J_{14}) ~,\\\\\n\\mathbb{W}^{03} = \\mathbb{W}^{43} = -EJ_{12}~, \\qquad & \\mathbb{W}^{31} = E(J_{20} +J_{24}) ~,\\\\\n\\mathbb{W}^{04} = & \\ 0 ~.\n\\end{split} \n\\end{equation}\nIf we rescale these generators and define:\n\\begin{equation}\\label{eq:eq7}\n\\begin{split}\n\\mathbb{R}_1 \\equiv \\frac{1}{E}\\mathbb{W}^{23}~, \\qquad \\mathbb{R}_{2} \\equiv & \\frac{1}{E}\\mathbb{W}^{31}~, \\qquad \\mathbb{R}_{3} \\equiv \\frac{1}{E}\\mathbb{W}^{12} ~,\\\\\n\\qquad \\cJ_{i} \\equiv &- \\frac{1}{E} \\mathbb{W}^{0i} ~,\n\\end{split}\n\\end{equation}\nthen these new operators satisfy:\n\\begin{equation}\n\\big[\\cJ_{i}, \\cJ_{j} \\big] = {\\rm i}\\varepsilon_{ijk}\\cJ_k~, \\qquad \n\\big[\\cJ_i , \\mathbb{R}_j \\big] = {\\rm i}\\varepsilon_{ijk} \\mathbb{R}_k~ , \n\\qquad \\big[\\mathbb{R}_i , \\mathbb{R}_j \\big] = 0~.\n\\end{equation}\nThese are the commutation relations for the three-dimensional Euclidean algebra, $\\mathfrak{iso}(3)$. The irreducible unitary representations of $\\mathfrak{iso}(3)$ are labelled by a continuous parameter $\\mu^2$, corresponding to the value the Casimir operator $\\mathbb{R}^{i} \\mathbb{R}_{i}$ takes. Since $\\mathbb{R}_{i}$ commute among themselves the operators can be simultaneously diagonalised, and the eigenvectors $\\ket{r_{i}}$ taken as a basis. However the only restriction on these is that $r_i r^i = \\mu^2$, which for non-zero $\\mu^2$ permits a continuous basis and is thus an infinite dimensional representation. Because we want only finite-dimensional representations, we must take:\n\\begin{equation} \\label{eq:eq6}\n\\mu^2 = 0 \\quad \\Longrightarrow \\quad \\mathbb{R}_{i} = 0 \\quad \\Longleftrightarrow \\quad J_{0i} = -J_{4i}~.\n\\end{equation}\nWe are therefore restricted to those representations in which the translation component is trivial, and so only the generators $\\cJ_i$ remain, which generate the algebra $\\mathfrak{so}(3)$. The algebra of the little group on massless representations is thus $\\mathfrak{so}(3) $ which is isomorphic to $\\mathfrak{su}(2)$. \nAs stated previously, the irreducible representations of $\\mathfrak{su}(2)$ are labelled by a non-negative (half) integer $s$ and have a single Casimir operator $\\cJ^i \\cJ_i$ which takes the value $s(s+1)\\mathbbm{1}$. This analysis leads to \\eqref{main}.\n\nThe spin value of a massless representation can still be found using a `spin' operator.\nThe following relation holds on massless representations:\n\\begin{equation}\n \\mathbb{S}_a :=- \\frac{1}{4} \\varepsilon_{abcde} J^{bc} \\mathbb{W}^{de} \n = \\cJ^2 P_{a} = s(s+1)P_a~,\n \\label{2.12}\n\\end{equation}\nwhere $\\cJ^2 = \\cJ^{i}\\cJ_{i}$ is the Casimir operator for the $\\mathfrak{so}(3)$ generators in \\eqref{eq:eq7}. The parameter $s$ is the spin of a massless particle.\nIts possible values in different representations are \n$s= 0, 1\/2, 1, $ and so on. Equation \\eqref{2.12} naturally holds for massless spinor and vector fields \\cite{Pindur}.\n\n\n\nIn general, the operator ${\\mathbb S}_a$ is not translationally invariant,\n\\begin{equation}\n\\big[ \\mathbb{S}_b, P_{a} \\big] = \\frac{{\\rm i}}{2}\\varepsilon_{abcde}P^{c}\\mathbb{W}^{de} ~.\n\\end{equation}\nIt is only for the massless representations with finite spin that the quantity on the right vanishes so that the spin operator commutes with the momentum operators. \nEquation \\eqref{2.12} is the five-dimensional analogue of the operator equation \\eqref{1.5}. Its consistency condition is \\eqref{main}.\\footnote{The consistency condition for \\eqref{1.5} is $P^{[a} {\\mathbb W}^{b]} =0$, which is the four-dimensional counterpart of \\eqref{main}.}\n\n\n\\section{Generalisations} \n\nThe results of section \\ref{section2.2} can be generalised to $d>5$ dimensions. \nThe Pauli-Lubanski tensor \\eqref{2.1} turns into\n\\begin{eqnarray}\n\\mathbb{W}^{a_1 \\dots a_{d-3}} = \\frac{1}{2}\\ve^{a_1 \\dots a_{d-3} bc e}J_{bc}P_e~.\n\\end{eqnarray}\nThe condition \\eqref{main} is replaced with\n\\begin{eqnarray}\nP^{[a} {\\mathbb W}^{b_1 \\dots b_{d-3}]} =0~.\n\\label{maind}\n\\end{eqnarray}\nThis equation is very similar to another that has appeared in the literature\nusing the considerations of conformal invariance \n\\cite{Bracken,BrackenJ,Siegel,SiegelZ}. One readily checks that \\eqref{maind} is equivalent to \n\\begin{eqnarray}\nJ_{ab} P^2 + 2 J_{c[a} P_{b]} P^c =0 \\quad \\implies \\quad \nJ_{c[a} P_{b]} P^c =0 ~.\n\\label{maind2}\n\\end{eqnarray}\nThe latter is solved on the momentum eigenstates \nby $J_{ab} p^b \\propto p_a$, which is of the form considered\nin \\cite{Bracken,BrackenJ,Siegel,SiegelZ}.\\footnote{We are grateful to Warren Siegel for useful comments.}\n\n\nEquation \\eqref{maind} characterises all irreducible massless representations of \n$\\mathsf{ISO}_0(d-1,1)$ with a finite (discrete) spin. Finally, the spin equation \\eqref{2.12}\nturns into \n\\begin{eqnarray}\n \\mathbb{S}_a := \\frac{(-1)^d}{2 (d-3)!} \\varepsilon_{abc e_1 \\dots e_{d-3}} J^{bc} \\mathbb{W}^{e_1 \\dots e_{d-3}} \n = \\cJ^2 P_{a} ~,\n\\end{eqnarray}\nwhere $\\cJ^2 = \\frac12 \\cJ^{ij} \\cJ_{ij} $ is the quadratic Casimir operator of the algebra $\\mathfrak{so}(d-2)$, with $i,j=1, \\dots , d-2$.\nFor every irreducible massless representation of \n$\\mathsf{ISO}_0(d-1,1)$ with a finite spin, it holds that $\\cJ^2 \\propto {\\mathbbm 1}$. \n\nIt is possible to derive a five-dimensional analogue of the operator equation \ndefining the $\\cN=1$ superhelicity $\\kappa$ in four dimensions \\cite{BK}. \nThe latter has the form\\footnote{In the supersymmetric case, the conventions of \\cite{BK} are used, \nin particular the Levi-Civita tensor $\\ve_{abcd}$ is normalised by $\\ve_{0123}=-1$.} \n\\begin{eqnarray}\n{\\mathbb L}_a = \\Big( \\kappa +\\frac 14 \\Big) P_a~, \n\\label{3.4}\n\\end{eqnarray}\nwhere the operator ${\\mathbb L}_a$ is defined by \n\\begin{eqnarray}\n{\\mathbb L}_a = {\\mathbb W}_a - \\frac{1}{16} (\\tilde \\sigma_a)^{\\ad \\alpha} \\big[ Q_\\alpha , \\bar Q_\\ad\\big] ~.\n\\label{3.5}\n\\end{eqnarray}\nThe fundamental properties of the operator ${\\mathbb L}_a $\n(the latter differs from the supersymmetric Pauli-Lubanksi vector \\cite{SS})\nare that it is translationally invariant and commutes with the supercharges $Q_\\alpha$ and $\\bar Q_\\ad$ in the massless representations of the $\\cN=1$ super-Poincar\\'e group.\\footnote{The irreducible massless\n representation of superhelicity $\\kappa$ is the direct sum of two irreducible massless Poincar\\'e representations corresponding to the helicity values $\\kappa$ and $\\kappa+\\frac12$.} \nThe superhelicity operator \\eqref{3.5} was generalised to higher dimensions\nin \\cite{PZ,AMT}. Generalisations of \\eqref{3.4} to five and higher dimension will be discussed elsewhere.\n\\\\\n\n\\noindent\n{\\bf Acknowledgements:}\\\\\nWe thank Warren Siegel for pointing out important references, and Michael Ponds for comments on the manuscript. \nSMK is grateful to Ioseph Buchbinder for email correspondence, and to Arkady Segal for discussions. The work of SMK work is supported in part by the Australian \nResearch Council, project No. DP200101944.\n\n\n\\begin{footnotesize}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nRecent advances on information technologies facilitate real-time message exchanges and decision-making among geographically dispersed strategic entities. This has boosted the emergence of new generation of networked systems; e.g., the smart grid and intelligent transportation systems. These networked systems share some common features: on one hand, the entities do not belong to a single authority and may pursue different or even competitive interests; on the other hand, each entity keeps private information which is unaccessible to others. It is of great interest to design practical mechanisms which allow for efficient coordination of self-interested entities and ensure network-wide performance. Game theory along with its distributed computation algorithms represents a promising tool to achieve the goal.\n\nIn many applications, distributed computation is executed in uncertain environments. For example, mobile robots are deployed in an operating environment where environmental distribution functions are unknown to robots in advance; e.g.,~\\cite{Stankovic.Johansson.Stipanovic:11,Zhu.Martinez:SICON13}. In traffic pricing, pricing policies of system operators may not be available to drivers. In optimal power flow control, the structural parameters of power systems are of national security interest and kept confidential from the public. The absence of such information makes game components; e.g., objective and constraint functions, inaccessible to players. Very recently, the informational constraint has been stimulating the investigation of \\emph{adaptive} algorithms, including~\\cite{Frihauf.Krstic.Basar:11,SL-MK:11,Stankovic.Johansson.Stipanovic:11} for continuous games and~\\cite{JRM-HPY-GA-JSS:06,Zhu.Martinez:SICON13} for discrete games.\n\n\\emph{Literature review.} Non-cooperative game theory has been widely used as a mathematical framework to reason about multiple selfish decision makers; see for instance~\\cite{Basar.Olsder:82}. These games\nhave found a variety of applications in economics, communication and robotics; see~\\cite{Altman.Basar:98,Chung.Hollinger.Isler:11,Dockner.Jorgensen.Long.Sorger:06,Frihauf.Krstic.Basar:11,Mitchell.Bayen.Tomlin:05}. In non-cooperative games, decision making of individuals is inherently distributed. Very recently, this attractive feature has been utilized to synthesize cooperative control schemes, and a partial reference list for this regard includes~\\cite{Arsie.Savla.ea:TAC09,Arslan.Marden.Shamma:07a,Li.Marden:11,Stankovic.Johansson.Stipanovic:11,Zhu.Martinez:SICON13}.\n\nThe set of papers more relevant to our work is concerned with generalized Nash games where strategy spaces are continuous and the actions of players are coupled through utility and constraint functions. Generalized Nash games are first formulated in~\\cite{Arrow.Debreu:54}. Since then, a great effort has been dedicated to studying the existence and structural properties of generalized Nash equilibria in; e.g.,~\\cite{Rosen:65} and the recent survey paper~\\cite{Facchinei.Kanzow:07}. A number of algorithms have been proposed to compute generalized Nash equilibria, including ODE-based methods~\\cite{SL-TB:87,Rosen:65}, nonlinear Gauss-Seidel-type approaches~\\cite{Pang.Scutari.Facchinei.Wang:08}, iterative primal-dual Tikhonov schemes~\\cite{Yin.Shanbhag.Mehta:11} and best-response dynamics~\\cite{Palomar.Eldar:10}.\n\nAs mentioned, the set of papers~\\cite{Frihauf.Krstic.Basar:11,SL-MK:11,JRM-HPY-GA-JSS:06,Stankovic.Johansson.Stipanovic:11,Zhu.Martinez:SICON13}\ninvestigates the \\emph{adaptiveness} of game theoretic learning algorithms. However, none of the papers mentioned in the last two paragraphs studies the \\emph{robustness} of the algorithms with respect to network unreliability; e.g., data transmission delays, quantization and dynamically changing topologies. In contrast, the robustness has been extensively studied for consensus and distributed optimization, including, to name a few,~\\cite{Jadbabaie.Lin.Morse:03,Nedic.Ozdaglar.Parrilo:08} for time-varying topologies,~\\cite{Rabbat.Nowak:05} for quantization and~\\cite{Munz.Papachristodoulou.Allgower:10} for time delays. Yet the adaptiveness issue has not been addressed in this group of papers.\n\n\n\\emph{Contributions.} In this paper, we aim to solve a class of generalized convex games over unreliable networks where the structures of component functions are unknown to the associated players. That is, we aim to simultaneously address the issues of adaptiveness and robustness for generalized convex games.\n\nIn the games, each player is associated with a convex objective function and subject to a private convex inequality constraint and a private convex constraint set. The component functions are assumed to be smooth and are unknown to the associated players. We investigate distributed first-order gradient-based computation algorithms for the following two scenarios:\n\n[Scenario One] The game map is pseudo-monotone and the maximum delay (equivalently, the maximum number of packet dropouts or link breaks) is bounded but unknown;\n\n[Scenario Two] The inequality constraints are absent, the (reduced) game map is strongly monotone and the maximum delay is known.\n\nInspired by simultaneous perturbation stochastic approximation for optimization in~\\cite{JS:03}, we utilize finite differences with diminishing approximation errors to estimate first-order gradients. We propose two distributed algorithms for the two scenarios and formally prove their asymptotic convergence. The comparison of the two proposed algorithms is given in Section~\\ref{sec:comparison}. The analysis integrates the tools from convex analysis, variational inequalities and simultaneous perturbation stochastic approximation. The algorithm performance is verified through demand response on the IEEE 30-bus Test System. A preliminary version of the current paper was published in~\\cite{Zhu.Frazzoli:CDC12} where the adaptiveness issue was not investigated.\n\n\n\\section{Problem formulation}\\label{sec:formulation}\n\nIn this section, we present the generalized convex game considered in the paper. It is followed by the notions and notations used throughout the paper.\n\n\\subsection{Generalized convex game}\n\nConsider the set of players $V \\triangleq \\{1,\\cdots,N\\}$ where the state of player~$i$ is denoted as $x^{[i]}\\in X_i\\subseteq \\mathds{R}^{n_i}$. The players are selfish and pursue different interests. In particular, given the joint state $x^{[-i]}\\in X_{-i} \\triangleq \\prod_{j\\neq i}X_j$ of its rivals\\footnote{We use the shorthand $-i\\triangleq V\\setminus\\{i\\}$ throughout the paper.}, each player~$i$ aims to solve the following program parameterized by $x^{[-i]}\\in X_{-i}$: \\begin{align}\\min_{x^{[i]}\\in X_i}f_i(x^{[i]},x^{[-i]}),\\quad {\\rm s.t.} \\quad G^{[i]}(x^{[i]},x^{[-i]})\\leq0,\\label{e10}\\end{align} where $f_i : \\mathds{R}^n\\rightarrow \\mathds{R}$ and $G^{[i]} : \\mathds{R}^n\\rightarrow\\mathds{R}^{m_i}$ with $n\\triangleq \\sum_{i\\in V}n_i$. In the remainder of the paper, we assume that the following properties about problem~\\eqref{e10} hold:\n\n\\begin{assumption} The maps $f_i$ and $G^{[i]}$ are smooth, and the maps $f_i(\\cdot,x^{[-i]})$ and $G^{[i]}(\\cdot,x^{[-i]})$ are convex in $x^{[i]}$. The set $X_i$ is convex and compact, and $X\\cap Y\\neq\\emptyset$ where $X \\triangleq \\prod_{i\\in V}X_i$ and $Y \\triangleq \\prod_{i\\in V}Y_i$ with $Y_i \\triangleq \\{x\\in X\\; |\\; G^{[i]}(x) \\leq 0\\}$.\\label{asm7}\n\\end{assumption}\n\n\n\nWe now proceed to provide an equivalent form of problem~\\eqref{e10}. To achieve this, we define the set-valued map $X_i^f : X_{-i} \\rightarrow 2^{X_i}$ as follows: \\begin{align*}X_i^f(x^{[-i]}) = \\{x^{[i]}\\in X_i \\; | \\; G^{[i]}(x^{[i]},x^{[-i]})\\leq 0\\}.\\end{align*} The set $X_i^f(x^{[-i]})$ represents the collection of feasible actions for player~$i$ when its opponents choose the joint state of $x^{[-i]}\\in X_{-i}$. With the map $X_i^f$, problem~\\eqref{e10} of player~$i$ is equivalent to the following one: \\begin{align}\\min_{x^{[i]}\\in X_i^f(x^{[-i]})}f_i(x^{[i]},x^{[-i]}).\\label{e11}\\end{align}\n\nGiven $x^{[-i]}\\in X_{-i}$, each player~$i$ aims to solve problem~\\eqref{e11}. The collection of such coupled optimization problems consists of the \\emph{generalized convex game} (for short, CVX). For the CVX game, we adopt the \\emph{generalized Nash equilibrium} (for short, GNE) as the solution notion which none of the players is willing to unilaterally deviate from:\n\n\\begin{definition} The joint state $\\tilde{x}\\in X\\cap Y$ is a generalized Nash equilibrium of the CVX game if the following holds: \\begin{align}&f_i(\\tilde{x})\\leq f_i(x^{[i]},\\tilde{x}^{[-i]}),\\quad \\forall x^{[i]}\\in X_i^f(\\tilde{x}^{[-i]}),\\quad \\forall i\\in V.\\nonumber\\end{align}\\label{def1}\n\\end{definition}\n\nDenote by $\\mathbb{X}_{\\rm CVX}$ the set of GNEs of the CVX game. The following lemma verifies the non-emptiness of $\\mathbb{X}_{\\rm CVX}$.\n\n\\begin{lemma} The set of generalized Nash equilibria of the CVX game is not empty, i.e., $\\mathbb{X}_{\\rm CVX} \\neq \\emptyset$.\\label{lem3}\n\\end{lemma}\n\n\\textbf{Proof:} Recall that $f_i$ is convex and $X\\cap Y$ is compact. Hence, $\\mathbb{X}_{\\rm CVX} \\neq \\emptyset$ is a direct result of~\\cite{Facchinei.Kanzow:07,Rosen:65}. \\oprocend\n\n\nIn the CVX game, the players desire to seek a GNE. It is noted that $f_i$, $G^{[i]}$, and $X_i$ are private information of player~$i$ and unaccessible to others. In order to compute a GNE, it becomes necessary that the players are inter-connected and able to communicate with each other to exchange their partial estimates of GNEs. The interconnection between players will be represented by a directed graph ${\\mathcal{G}} = (V,\\mathcal{E})$ where\n$\\mathcal{E}\\subset V\\times V\\setminus {\\rm diag}(V)$ is the set of\nedges. The neighbor relation is determined by the dependency of $f_i$ and\/or $G^{[i]}$ on $x^{[j]}$. In particular, $(i,j)\\in\\mathcal{E}$ if and only if $f_i$ and\/or $G^{[i]}$ depend upon $x^{[j]}$. Denote by $\\mathcal{N}_i^{\\rm IN} \\triangleq \\{j\\in V\\;|\\; (i,j)\\in\\mathcal{E}\\}$ the set of in-neighbors of\nplayer~$i$.\n\nIn this paper, we aim to develop distributed algorithms which allow for the computation of GNEs in the presence of the following two challenges.\n\\begin{enumerate}\n\\item Data transmissions between the players in $V$ are unreliable, and subject to transmission delays, packet dropouts and\/or link breaks. We let $x^{[j]}(k - \\tau^{[i]}_j(k))$ with $\\tau^{[i]}_j(k)\\geq0$ be either $(i)$ the outdated state of player~$j$ received by player~$i$ at time~$k$ due to transmission delays; or $(ii)$ the latest state of player~$j$ received by player~$i$ by time~$k$ due to packet dropouts and\/or link breaks. For example, if the link from player $j$ to player $i$ is broken at time $k$, then player~$i$ cannot receive the message $x^{[j]}(k)$ from player $j$. For this case, player~$i$ uses the more recent received message sent from player $j$ which is $x^{[j]}(k-\\tau_j^{[i]}(k))$. Denote by $\\Lambda_i(k)\\triangleq \\{x^{[j]}(k-\\tau^{[i]}_j(k))\\}_{j\\in\\mathcal{N}^{\\rm IN}_i}$ the set of latest states of $\\mathcal{N}_i^{\\rm IN}$ received by player~$i$ at time~$k$. Let $\\tau_{\\max} \\triangleq \\sup_{k\\geq0}\\max_{i\\in V}\\max_{j\\in\\mathcal{N}_i^{\\rm IN}}\\tau^{[i]}_j(k)$ be the maximum delay or the maximum number of successive packet dropouts.\n\\item Each player~$i$ is unaware of the structures of $f_i$ and $G^{[i]}$ but can observe their realizations. That is, if the players input the joint state $x$ into $f_i$, then player~$i$ can observe the realized value of $f_i(x)$.\n\\end{enumerate}\n\nThe disclosure of the value $f_i(x)$ is case dependent. In~\\cite{Stankovic.Johansson.Stipanovic:11,Zhu.Martinez:SICON13}, mobile sensors are unaware of environmental distribution functions but they can observe induced utilities via on-site measurements. This is an example of engineering systems. Section~\\ref{sec:simulations} will provide a concrete example of demand response of power networks where the system operator discloses realized values via communication given the decisions of end-users. This is an example of social systems.\n\n\n\n\n\n\n\n\n\\section{Preliminaries}\\label{sec:preliminaries}\n\nIn the CVX game, each player is subject to an inequality constraint. In optimization literature, Lagrangian relaxation is a widely used approach to handle inequality constraints. Following this vein, we will perform Lagrangian relaxation on the CVX game and obtain the unconstrained convex (for short, UC) game.\n\nIn the UC game, there are two sets of players: the set of primal players $V$ and the set of dual players $V_m \\triangleq \\{1,\\cdots,N\\}$. Define the following private Lagrangian $\\mathcal{H}_i : \\mathds{R}^n\\times\\mathds{R}^{m_i}_{\\geq0}\\rightarrow\\mathds{R}$: $\\mathcal{H}_i(x,\\mu^{[i]}) \\triangleq f_i(x) + \\langle \\mu^{[i]}, G^{[i]}(x)\\rangle$, for primal player~$i\\in V$ and dual player~$i\\in V_m$. Given any $x^{[-i]}\\in X_{-i}$ and $\\mu^{[i]}\\in\\mathds{R}^{m_i}_{\\geq0}$, each primal player~$i\\in V$ aims to minimize $\\mathcal{H}_i$ over $x^{[i]}\\in X_i$; i.e., $\\min_{x^{[i]}\\in X_i}\\mathcal{H}_i(x^{[i]},x^{[-i]},\\mu^{[i]})$, and, instead, the objective of the dual player $i\\in V_m$ is to maximize $\\mathcal{H}_i$ over $\\mu^{[i]}\\in M_i\\subseteq\\mathds{R}^{m_i}_{\\geq0}$; i.e., $\\min_{\\mu^{[i]}\\in M_i}-\\mathcal{H}_i(x,\\mu^{[i]})$ where the set $M_i\\subseteq\\mathds{R}^{m_i}_{\\geq0}$ is convex, non-empty and will be introduced in the sequel. We let $\\eta \\triangleq (x,\\mu)$ and the set $K \\triangleq X\\times M$ with $M \\triangleq \\prod_{i\\in V}M_i$. The above game among the players in $V\\cup V_m$ is referred to as the UC game parameterized by the set $M$. The set of $M$ will play an important role in determining the properties of the UC game and we will discuss the choice of $M$ later. The solution concept for the UC game parameterized by $M$ is the standard notion of Nash equilibrium given below:\n\n\\begin{definition} The joint state of $(\\tilde{x},\\tilde{\\mu}) \\in X\\times M$ is a Nash equilibrium of the UC game parameterized by $M$ if the following holds for all $i\\in V$: $\\mathcal{H}_i(\\tilde{x},\\tilde{\\mu}^{[i]})\\leq \\mathcal{H}_i(x^{[i]},\\tilde{x}^{[-i]},\\tilde{\\mu}^{[i]}),\\quad \\forall x^{[i]}\\in X_i$, and the following holds for all $i\\in V_m$: $\\mathcal{H}_i(\\tilde{x},\\mu^{[i]})\\leq \\mathcal{H}_i(\\tilde{x},\\tilde{\\mu}^{[i]}),\\quad \\forall \\mu^{[i]}\\in M_i$.\\label{def2}\n\\end{definition} Denote by $\\mathbb{X}_{\\rm UC}(M)$ the set of NEs of the UC game parameterized by $M$. We now proceed to illustrate the relation between the UC game and the CVX game. Before doing so, let us state the following boundedness assumption on the dual solutions:\n\n\\begin{assumption} There is a vector $\\vartheta \\triangleq (\\vartheta_i)_{i\\in V}\\in\\mathds{R}^N_{>0}$ such that for any $(\\tilde{x},\\tilde{\\mu})\\in \\mathbb{X}_{\\rm UC}(\\mathds{R}^m_{\\geq0})$, $\\|\\tilde{\\mu}^{[i]}\\| \\leq \\vartheta_i$ for $i\\in V$.\\label{asm1}\n\\end{assumption}\n\n\\begin{remark} We would like to make a remark on Assumption~\\ref{asm1}. For convex optimization, it is shown in~\\cite{Hiriart.Lemarechal:96} that the Lagrangian multipliers are uniformly bounded under the standard Slater's condition. This boundedness property is used in~\\cite{Nedic.Ozdagalr:08} and further in~\\cite{Zhu.Martinez:09tac,Zhu.Martinez:10tac} to solve convex and non-convex programs in centralized and distributed manners. In this paper, Assumption~\\ref{asm1}, on one hand, ensures the existence of GNEs, and on the other hand, guarantees the boundedness of estimates and gradients for the case of unknown $\\tau_{\\max}$. We will discuss the verification of Assumption~\\ref{asm1} in Section~\\ref{sec:discussion}.\\oprocend\\label{rem6}\n\\end{remark}\n\nThe following proposition characterizes the relations between the UC game and the CVX game.\n\n\\begin{proposition} The following properties hold:\n\\begin{enumerate}\n\\item[(P1)] Consider any $(\\tilde{x}, \\tilde{\\mu})\\in\\mathbb{X}_{\\rm UC}(M)$. We have $\\tilde{x}\\in \\mathbb{X}_{\\rm CVX}$ if the following properties hold:\n\\begin{itemize}\n\\item feasibility; i.e., $G^{[i]}(\\tilde{x})\\leq0$ for all $i\\in V$;\n\\item slackness complementarity; i.e., $\\langle \\tilde{\\mu}^{[i]}, G^{[i]}(\\tilde{x})\\rangle = 0$ for all $i\\in V$.\n\\end{itemize}\n\\item[(P2)] Suppose Assumption~\\ref{asm1} holds. Consider the UC game parameterized by $M$ with $M_i \\triangleq \\{\\mu^{[i]}\\in\\mathds{R}^{m_i}_{\\geq0}\\;|\\;\\|\\mu^{[i]}\\|\\leq \\vartheta_i+r_i\\}$ for some $r_i>0$. Then $\\mathbb{X}_{\\rm UC}(M)\\neq\\emptyset$. In addition, for any $(\\tilde{x}, \\tilde{\\mu})\\in\\mathbb{X}_{\\rm UC}(M)$, it holds that $\\tilde{x}\\in \\mathbb{X}_{\\rm CVX}$.\n\\end{enumerate} \\label{pro5}\n\\end{proposition}\n\n\\textbf{Proof:} The proof of (P1) is a slight extension of the results in~\\cite{Bertsekas:09} to our game setup. For the sake of completeness, we summarize the analysis here. Since $\\langle \\tilde{\\mu}^{[i]}, G^{[i]}(\\tilde{x})\\rangle = 0$, we have the following relation: \\begin{align}f_i(\\tilde{x}) = \\mathcal{H}_i(\\tilde{x},\\tilde{\\mu}^{[i]}) - \\langle\\tilde{\\mu}^{[i]},G^{[i]}(\\tilde{x})\\rangle = \\mathcal{H}_i(\\tilde{x},\\tilde{\\mu}^{[i]}).\\label{e7}\\end{align} Since $(\\tilde{x},\\tilde{\\mu})\\in\\mathbb{X}_{\\rm UC}(M)$, then the following relation holds for all $x^{[i]}\\in X_i$: \\begin{align}\\mathcal{H}_i(\\tilde{x},\\tilde{\\mu}^{[i]})\n-\\mathcal{H}_i(x^{[i]},\\tilde{x}^{[-i]},\\tilde{\\mu}^{[i]}) \\leq 0.\\label{e17}\\end{align} Substitute~\\eqref{e7} into~\\eqref{e17}, and it renders the following for any $x^{[i]}\\in X_i$:\n\\begin{align}f_i(\\tilde{x}) \\leq f_i(x^{[i]},\\tilde{x}^{[-i]}) + \\langle\\tilde{\\mu}^{[i]},G^{[i]}(x^{[i]},\\tilde{x}^{[-i]})\\rangle,\\nonumber\\end{align} and thus $f_i(\\tilde{x}) \\leq f_i(x^{[i]},\\tilde{x}^{[-i]}),\\quad \\forall x^{[i]}\\in X_i^f(\\tilde{x}^{[-i]})$, by noting that $\\langle\\tilde{\\mu}^{[i]},G^{[i]}(x^{[i]},\\tilde{x}^{[-i]})\\rangle \\leq 0$ for any $ x^{[i]}\\in X_i^f(\\tilde{x}^{[-i]})$. Recall that $G^{[i]}(\\tilde{x})\\leq0$ for all $i\\in V$. The above arguments hold for all $i\\in V$, and thus it establishes that $\\tilde{x}\\in\\mathbb{X}_{\\rm CVX}$.\n\nWe now proceed to show (P2). Since $\\mu^{[i]}\\geq0$, $\\mathcal{H}_i$ is a positive combination of convex functions of $x^{[i]}$. Thus $\\mathcal{H}_i$ is convex in $x^{[i]}$. It is easy to see that $-\\mathcal{H}_i$ is convex (actually affine) in $x^{[i]}$. Since $X$ and $M$ are convex and compact, it follows that $\\mathbb{X}_{\\rm UC}(M) \\neq \\emptyset$ by the results on generalized convex games in; e.g.,~\\cite{Facchinei.Kanzow:07,Rosen:65}. Pick any $(\\tilde{x},\\tilde{\\mu})\\in\\mathbb{X}_{\\rm UC}(M)$ and then $\\|\\tilde{\\mu}^{[i]}\\| \\leq \\vartheta_i$ by Assumption~\\ref{asm1}. We then have the following relation: \\begin{align}{\\mathcal{H}}_i(\\tilde{x},\\tilde{\\mu}^{[i]})\n -{\\mathcal{H}}_i(\\tilde{x},\\mu^{[i]}) \\geq 0,\\quad \\forall\\mu^{[i]}\\in M_i,\\nonumber\\end{align} and thus,\n\\begin{align}\\langle \\mu^{[i]}-\\tilde{\\mu}^{[i]}, G^{[i]}(\\tilde{x})\\rangle \\leq 0,\\quad \\forall\\mu^{[i]}\\in M_i.\\label{e6}\\end{align}\n\nWe now proceed to verify by contradiction the feasibility of $G^{[i]}(\\tilde{x}) \\leq 0$ and the complementary slackness of $\\langle \\tilde{\\mu}^{[i]}, G^{[i]}(\\tilde{x})\\rangle = 0$. Assume that $G^{[i]}(\\tilde{x})_{\\ell} > 0$. We choose $\\mu^{[i]}$ such that $\\mu^{[i]}_{\\ell} = \\tilde{\\mu}^{[i]}_{\\ell} + \\pi r_i$ and $\\mu^{[i]}_{\\ell'} = 0$ for $\\ell' \\neq \\ell$ where $\\pi>0$ is sufficiently small such that $\\mu^{[i]}\\in M_i$. Then it follows from~\\eqref{e6} that $r_i G^{[i]}(\\tilde{x})_{\\ell} \\leq 0$ which is a contradiction. Hence, we have the feasibility of $G^{[i]}(\\tilde{x})\\leq 0$. Combine it with $\\tilde{\\mu}^{[i]}\\geq0$, and it renders that $\\langle \\tilde{\\mu}^{[i]}\\rangle, G^{[i]}(\\tilde{x})\\leq 0$. On the other hand, we let $\\mu^{[i]} = 0$ in the relation~\\eqref{e6} and have $\\langle \\tilde{\\mu}^{[i]}, G^{[i]}(\\tilde{x})\\rangle \\geq 0$. Hence, it renders that $\\langle \\tilde{\\mu}^{[i]}, G^{[i]}(\\tilde{x})\\rangle = 0$. We reach the desired result by using the property (P1). \\oprocend\n\n\\subsection{Notations and notions}\\label{subsection:notations}\n\nThe vector $\\textbf{1}_n$ represents the column vector with $n$ ones. The vector $e^{[\\ell]}_{n}$ is the one in $\\mathds{R}^n$ whose $\\ell$-th coordinate is one and whose other coordinates are all zero. For any pair of vectors $a,b\\in\\mathds{R}^p$, the relation $a\\leq b$ means $a_{\\ell}\\leq b_{\\ell}$ for all $1\\leq\\ell\\leq p$. Since the function $f_i$ is continuous and $X_i$ is compact, then the following quantities are well-defined:\n$\\sigma_{i,\\min} = \\inf_{x\\in X_i}f_i(x),\\quad \\sigma_{i,\\max} = \\sup_{x\\in X_i}f_i(x),\\quad\n\\sigma_{\\min} = \\inf_{x\\in X}\\sum_{i\\in V}f_i(x),\\quad \\sigma_{\\max} = \\sup_{x\\in X}\\sum_{i\\in V}f_i(x)$. Given a non-negative scalar sequence $\\{\\alpha(k)\\}_{k\\geq0}$, it is \\emph{summable} if $\\sum_{k=0}^{\\infty}\\alpha(k)<+\\infty$ and \\emph{square summable} if $\\sum_{k=0}^{\\infty}\\alpha(k)^2<+\\infty$.\n\nIn the remainder of the paper, we will use some notions about \\emph{monotonicity}, and the readers are referred to~\\cite{Aubin.Frankowska:09,Facchinei.Pang:03} for detailed discussion. The mapping $F : Z \\rightarrow Z'$ is \\emph{strongly monotone} with constant $\\rho > 0$ over $Z$ if for each pair of $\\eta,\\eta'\\in Z$, the following holds: \\begin{align*}\\langle F(\\eta) - F(\\eta'), \\eta - \\eta' \\rangle \\geq \\rho \\|\\eta - \\eta'\\|^2.\\end{align*} The mapping $F : Z \\rightarrow Z'$ is \\emph{monotone} over $Z$ if for each pair of $\\eta,\\eta'\\in Z$, the following holds: \\begin{align*}\\langle F(\\eta) - F(\\eta'), \\eta - \\eta' \\rangle \\geq 0.\\end{align*} The mapping $F : Z \\rightarrow Z'$ is \\emph{pseudo-monotone} over $Z$ if for each pair of $\\eta,\\eta'\\in Z$, it holds that $\\langle F(\\eta'), \\eta - \\eta' \\rangle \\geq 0$ implies $\\langle F(\\eta), \\eta - \\eta' \\rangle \\geq 0$. It is known that strong monotonicity implies monotonicity, and monotonicity implies pseudo-monotonicity~\\cite{Facchinei.Pang:03}.\n\nGiven the non-empty, convex, and closed set $Z\\in {\\mathds{R}}^n$, the\nprojection operator onto $Z$, $P_Z: \\mathds{R}^n \\rightarrow Z$, is\ndefined as $P_Z[z] = {\\rm argmin}_{x\\in Z}\\|x-z\\|$. The following is the non-expansiveness of the projection operator; e.g.,~\\cite{Bertsekas:09}.\n\n\\begin{lemma} Let $Z$ be a non-empty, closed and convex set in ${\\mathds{R}}^n$. For any $z\\in{\\mathds{R}}^n$, the following holds for any $y\\in Z$: $\\|P_X[z] - y\n \\|^2 \\leq \\| z - y \\|^2 - \\| P_X[z] - z\\|^2$. \\label{lem6}\n\\end{lemma}\n\nWe define $\\nabla_{x^{[i]}}\\mathcal{H}_i(x,\\mu^{[i]})$ as the gradient of the convex function $\\mathcal{H}_i(\\cdot,x^{[-i]},\\mu^{[i]})$ at $x^{[i]}$, and $\\nabla_{\\mu^{[i]}}\\mathcal{H}_i(x,\\mu^{[i]}) = G^{[i]}(x)$ as the gradient of the concave (actually affine) function $\\mathcal{H}_i(x,\\cdot)$ at $\\mu^{[i]}$. Define $\\nabla \\Theta$ as the map of partial gradients of the players' objective functions: \\begin{align*}\\nabla \\Theta(\\eta) &\\triangleq [\\nabla_{x^{[1]}}\\mathcal{H}_1(x,\\mu^{[1]})^T\\cdots \\nabla_{x^{[N]}}\\mathcal{H}_N(x,\\mu^{[N]})^T\\\\ &\\nabla_{\\mu^{[1]}}-\\mathcal{H}_1(x,\\mu^{[1]})\\cdots \\nabla_{\\mu^{[N]}}-\\mathcal{H}_N(x,\\mu^{[N]})]^T.\\end{align*} The map $\\nabla \\Theta$ is referred to as the \\emph{game map}. It is well-known that $\\nabla\\Theta$ is monotone (resp. strongly monotone) over $K$ if and only if its Jacobian $\\nabla^2\\Theta(\\eta)$ is positive semi-definite (resp. positive definite) for all $\\eta\\in K$.\n\nFor $G^{[i]}(x)\\leq0$ being absent, we define $\\nabla_{x^{[i]}}f_i(x)$ as the gradient of the convex function $f_i(\\cdot,x^{[-i]})$ at $x^{[i]}$. Define $\\nabla \\Theta^r$ as the map of partial gradients of the players' objective functions: \\begin{align*}\\nabla \\Theta^r(x) \\triangleq [\\nabla_{x^{[1]}}f_1(x)^T\\cdots \\nabla_{x^{[N]}}f_N(x)^T]^T.\\end{align*} The map $\\nabla \\Theta^r$ is referred to as the \\emph{(reduced) game map} for $G^{[i]}(x)\\leq0$ being absent.\n\n\nGiven any convex and closed set $Z$, let $\\mathbb{P}_Z$ be the projection operator onto the set $Z$.\n\nThe following lemma is a direct result of the smoothness of the component functions and the compactness of the constraint sets.\n\n\\begin{lemma} The image set of the game map $\\nabla\\Theta$ is uniformly bounded over $K$. In addition, the game map $\\nabla\\Theta$ (resp. $\\nabla\\Theta^r$) is Lipschitz continuous over $K$ with constant $L_{\\Theta}$ (resp. $L_{\\Theta^r}$). \\label{lem4}\n\\end{lemma}\n\n\\section{Distributed computation algorithms}\\label{sec:algorithm}\n\nIn this section, we will study the scenarios mentioned in the introduction. Proposition~\\ref{pro5} reveals that it suffices to compute a Nash equilibrium in $\\mathbb{X}_{\\rm UC}$ if Assumption~\\ref{asm1} holds. In the remainder of this section, we will suppose Assumption~\\ref{asm1} holds, and then each dual player~$i$ can define the set of $M_i \\triangleq \\{\\mu^{[i]}\\in\\mathds{R}^{m_i}_{\\geq0}\\;|\\;\\|\\mu^{[i]}\\|\\leq \\vartheta_i+r_i\\}$ with $r_i>0$. Assumption~\\ref{asm1} will be justified in Section~\\ref{sec:assumption}.\n\n\\subsection{Scenario one: pseudo-monotone game map and unknown $\\tau_{\\max}$}\n\nIn this section, we synthesize a distributed first-order algorithm for the case where the quantity $\\tau_{\\max}$ is unknown and the game map $\\nabla \\Theta$ is merely pseudo-monotone.\n\nAlgorithm~\\ref{ta:algo1} is based on projected primal-dual gradient methods where each primal or dual player~$i$\\footnote{In practical implementation, the update rules of primal player~$i$ and dual player~$i$ are both executed by real player~$i$.} updates its own estimate by moving along its partial gradient with a certain step-size and projecting the estimate onto the local constraint set. Recall that $f_i$ and $G^{[i]}$ are unknown to player~$i$. Then player~$i$ cannot compute partial gradient $\\nabla_{x^{[i]}}\\mathcal{H}_i$. Inspired by simultaneous perturbation stochastic approximation in; e.g.,~\\cite{JS:03}, we use finite differences to approximate the partial gradient $\\nabla_{x^{[i]}}\\mathcal{H}_i$; i.e., for $\\ell = 1,\\cdots,n_i$, \\begin{align}&[\\nabla_{x^{[i]}}\\mathcal{H}_i(x(k),\\mu^{[i]}(k))]_{\\ell}\\nonumber\\\\\n&\\approx \\frac{1}{2c_i(k)}\\big(\\mathcal{H}_i(x^{[i]}(k)+c_i(k)e^{[\\ell]}_{n_i},x^{[-i]}(k),\\mu^{[i]}(k))\\nonumber\\\\\n&-\\mathcal{H}_i(x^{[i]}(k)-c_i(k)e^{[\\ell]}_{n_i},x^{[-i]}(k),\\mu^{[i]}(k))\\big),\\label{e51}\\end{align}\nwhere $-c_i(k)e^{[\\ell]}_{n_i}$ and $c_i(k)e^{[\\ell]}_{n_i}$ are two-way perturbations. In addition, $\\nabla_{\\mu^{[i]}}\\mathcal{H}_i(x(k),\\mu^{[i]}(k)) = G^{[i]}(x(k))$. In Algorithm~\\ref{ta:algo1}, $\\mathcal{D}^{[i]}_x(k)$ (resp. $\\mathcal{D}^{[i]}_{\\mu}(k)$) is the right-hand side of~\\eqref{e51} (resp. $\\nabla_{\\mu^{[i]}}\\mathcal{H}_i(x,\\mu^{[i]})$) with delayed estimates. Here we assume that player~$i$ can observe the \\emph{values} associated with $f_i$ and $G^{[i]}$ and thus $\\mathcal{D}^{[i]}_x(k)$ and $\\mathcal{D}^{[i]}_{\\mu}(k)$.\n\nThe scalar $c_i(k)>0$ is the perturbation magnitude. When it is small, the finite-difference approximation is close to the partial gradient. In order to asymptotically eliminate the error induced by the finite-difference approximation, $c_i(k)$ needs to be diminishing. The decreasing rate of $c_i(k)$ should match that of computation step-sizes $\\alpha(k)$. Otherwise, the convergence to Nash equilibrium may be prevented. This is captured by (A4) of Theorem~\\ref{the1}.\n\n\\begin{algorithm}[htbp]\\caption{Distributed gradient-based algorithm for Scenario one} \\label{ta:algo1}\n\\textbf{Require:} Each primal player $i\\in V$ chooses the initial state $x^{[i]}(0)\\in X_i$. And each dual player~$i\\in V_m$ defines the set $M_i$ and chooses the initial state $\\mu^{[i]}(0)\\in M_i$.\n\n\\textbf{Ensure:} At each $k \\ge 0$, each player in $V\\cup V_m$ executes the following steps:\n\n1. Each primal player $i\\in V$ updates its state according to the following rule:\\begin{align*} x^{[i]}(k+1) = \\mathbb{P}_{X_i}[x^{[i]}(k) - \\alpha(k)\\mathcal{D}^{[i]}_x(k)],\n\\end{align*} where the approximate gradient $\\mathcal{D}^{[i]}_x(k)$ is given by \\begin{align}&[\\mathcal{D}^{[i]}_x(k)]_{\\ell} = \\frac{1}{2c_i(k)}\n\\{f_i(x^{[i]}(k)+c_i(k)e^{[\\ell]}_{n_i},\\Lambda_i(k))\\nonumber\\\\\n&+ \\langle\\mu^{[i]}(k),G^{[i]}(x^{[i]}(k)+c_i(k)e^{[\\ell]}_{n_i},\\Lambda_i(k))\\rangle\\nonumber\\\\\n&-f_i(x^{[i]}(k)-c_i(k)e^{[\\ell]}_{n_i},\\Lambda_i(k))\\nonumber\\\\\n&- \\langle\\mu^{[i]}(k),G^{[i]}(x^{[i]}(k)\n-c_i(k)e^{[\\ell]}_{n_i},\\Lambda_i(k))\\rangle\\},\\label{e38}\\end{align} for $\\ell = 1,\\cdots,n_i$.\n\n2. Each dual player $i\\in V_m$ updates its state according to the following rule:\\begin{align*}\n\\mu^{[i]}(k+1) = \\mathbb{P}_{M_i}[\\mu^{[i]}(k) + \\alpha(k)\\mathcal{D}^{[i]}_{\\mu}(k)],\n\\end{align*} where the gradient $\\mathcal{D}^{[i]}_{\\mu}(k)$ is given by $\\mathcal{D}^{[i]}_{\\mu}(k) = \\nabla_{\\mu^{[i]}}\\mathcal{H}_i(x^{[i]}(k),\\Lambda_i(k),\\mu^{[i]}(k)) = G^{[i]}(x^{[i]}(k),\\Lambda_i(k))$.\n\n3. Repeat for $k = k+1$.\n\\end{algorithm}\n\nThe following theorem demonstrates that Algorithm~\\ref{ta:algo1} is able to achieve a GNE from any initial state in $K$.\n\n\\begin{theorem} Suppose the following hold:\n\\begin{enumerate}\n\\item[(A1)] the quantity $\\tau_{\\max}$ is finite;\n\\item[(A2)] Assumptions~\\ref{asm7} and~\\ref{asm1} hold;\n\\item[(A3)] the sequence of $\\{\\alpha(k)\\}$ is positive, not summable but square summable;\n\\item[(A4)] the sequence of $\\displaystyle{\\{\\alpha(k)\\max_{i\\in V}c_i(k)\\}}$ is summable;\n\\item[(A5)] the game map $\\nabla\\Theta$ is pseudo-monotone over $K$.\n\\end{enumerate}\nFor any initial state $\\eta(0)\\in K$, the sequence of $\\{\\eta(k)\\}$ generated by Algorithm~\\ref{ta:algo1} converges to some $(\\tilde{x},\\tilde{\\mu})\\in{\\mathbb{X}}_{\\rm UC}$ where $\\tilde{x}\\in\\mathbb{X}_{\\rm CVX}$.\\label{the1}\n\\end{theorem}\n\n\\textbf{Proof:} First of all, it is noted that $K$ is compact since each $X_i$ and $X_i$ is compact in Assumption~\\ref{asm1} and (P2) in Proposition~\\ref{pro5} as well as $\\nabla\\Theta$ is uniformly bounded over $K$ in Lemma~\\ref{lem4}.\n\nThen we write Algorithm~\\ref{ta:algo1} in the following compact form:\n\\begin{align}\\eta(k+1) = \\mathbb{P}_K[\\eta(k)-\\alpha(k)\\mathcal{D}(k)],\\label{e19}\n\\end{align} where $\\mathcal{D}(k) \\triangleq ((\\mathcal{D}^{[i]}_x(k)^T)^T_{i\\in V},(-\\mathcal{D}^{[i]}_{\\mu}(k)^T)^T_{i\\in V_m})^T$ is subject to time delays and perturbations. Then pick any $\\eta\\in K$. It follows from the non-expansiveness of the projection operator $\\mathbb{P}_K$, Lemma~\\ref{lem6}, that for any $\\eta\\in K$, the following relation holds:\n\\begin{align}&\\|\\eta(k+1)-\\eta\\|^2\\nonumber\\\\\n&\\leq \\|\\eta(k)-\\alpha(k)\\mathcal{D}(k)-\\eta\\|^2-\\|\\mathcal{D}(k)\\|^2\\nonumber\\\\\n&\\leq \\|\\eta(k)-\\alpha(k)\\mathcal{D}(k)-\\eta\\|^2\\nonumber\\\\\n&= \\|\\eta(k)-\\eta\\|^2 - 2\\alpha(k)\\langle \\mathcal{D}(k), \\eta(k)-\\eta\\rangle\\nonumber\\\\\n&+ \\alpha(k)^2\\|\\mathcal{D}(k)\\|^2.\\label{e25}\n\\end{align}\n\nIt follows from~\\eqref{e25} that\n\\begin{align}&2\\alpha(k)\\langle \\mathcal{D}(k), \\eta(k)-\\eta\\rangle\\leq \\|\\eta(k) - \\eta\\|^2\\nonumber\\\\\n&- \\|\\eta(k+1) - \\eta\\|^2 +\\alpha(k)^2\\|\\mathcal{D}(k)\\|^2.\\label{e23}\\end{align}\n\nFor any $k\\geq0$ and $i$, we define the following: \\begin{align}&\\hat{\\mathcal{D}}^{[i]}_x(k) = \\nabla_{x^{[i]}}\\mathcal{H}_i(x^{[i]}(k),\\{x^{[j]}(k)\\}_{j\\in\\mathcal{N}^{\\rm IN}_i},\\mu^{[i]}(k)),\\nonumber\\\\\n&\\hat{\\mathcal{D}}^{[i]}_{\\mu}(k) \\triangleq \\nabla_{\\mu^{[i]}}\\mathcal{H}_i(x^{[i]}(k),\\{x^{[j]}(k)\\}_{j\\in\\mathcal{N}^{\\rm IN}_i},\\mu^{[i]}(k)),\\nonumber\\\\\n&\\hat{\\mathcal{D}}(k) \\triangleq ((\\hat{\\mathcal{D}}^{[i]}_x(k)^T)^T_{i\\in V},(-\\hat{\\mathcal{D}}^{[i]}_{\\mu}(k)^T)^T_{i\\in V_m})^T,\\label{e16}\\end{align} which are the gradients evaluated at the delay-free states. So, the quantity $\\hat{\\mathcal{D}}(k)$ is free of time delays and perturbations.\n\nSimilarly, for any $k\\geq0$ and $i$, we define the following: \\begin{align}&\\tilde{\\mathcal{D}}^{[i]}_x(k) = \\nabla_{x^{[i]}}\\mathcal{H}_i(x^{[i]}(k),\\Lambda_i(k),\\mu^{[i]}(k)),\\nonumber\\\\\n&\\tilde{\\mathcal{D}}^{[i]}_{\\mu}(k) \\triangleq \\nabla_{\\mu^{[i]}}\\mathcal{H}_i(x^{[i]}(k),\\Lambda_i(k),\\mu^{[i]}(k)),\\nonumber\\\\\n&\\tilde{\\mathcal{D}}(k) \\triangleq ((\\tilde{\\mathcal{D}}^{[i]}_x(k)^T)^T_{i\\in V},(-\\tilde{\\mathcal{D}}^{[i]}_{\\mu}(k)^T)^T_{i\\in V_m})^T,\\nonumber\\end{align} which are the gradients evaluated at the delayed states. Then the quantity $\\tilde{\\mathcal{D}}(k)$ is free of perturbations but subject to time delays.\n\nWith the above notations at hand, the relation~\\eqref{e23} implies the following: \\begin{align}&2\\alpha(k)\\langle \\hat{\\mathcal{D}}(k), \\eta(k)-\\eta\\rangle \\leq \\alpha(k)^2\\|\\mathcal{D}(k)\\|^2 + \\|\\eta(k) - \\eta\\|^2\\nonumber\\\\&- \\|\\eta(k+1) - \\eta\\|^2 + 2\\alpha(k)\\langle\\tilde{\\mathcal{D}}(k)-\\mathcal{D}(k),\\eta(k)-\\eta\\rangle\\nonumber\\\\\n&+ 2\\alpha(k)\\langle\\hat{\\mathcal{D}}(k)-\\tilde{\\mathcal{D}}(k),\\eta(k)-\\eta\\rangle\\nonumber\\\\\n&\\leq \\alpha(k)^2\\|\\mathcal{D}(k)\\|^2 + \\|\\eta(k) - \\eta\\|^2\\nonumber\\\\&- \\|\\eta(k+1) - \\eta\\|^2 + 2\\alpha(k)\\|\\tilde{\\mathcal{D}}(k)-\\mathcal{D}(k)\\| \\|\\eta(k)-\\eta\\|\\nonumber\\\\\n&+ 2\\alpha(k)\\|\\hat{\\mathcal{D}}(k)-\\tilde{\\mathcal{D}}(k)\\| \\|\\eta(k)-\\eta\\|.\\label{e30}\\end{align}\n\nNow let us examine the term with $\\|\\hat{\\mathcal{D}}(k)-\\tilde{\\mathcal{D}}(k)\\|$ on the right-hand side of~\\eqref{e30}. Since $x(k)\\in X$ and $X$ is convex and closed, it then follows from the non-expansiveness of the projection operator $\\mathbb{P}_X$, that \\begin{align}\\|\\eta(k+1)-\\eta(k)\\|\n&=\\|\\mathbb{P}_K[\\eta(k)-\\alpha(k)\\mathcal{D}(k)]-\\mathbb{P}_K[\\eta(k)]\\|\\nonumber\\\\\n&\\leq\\alpha(k)\\|\\mathcal{D}(k)\\|.\\label{e14}\\end{align} Consequently,\nit follows from the Lipschitiz continuity of the game map $\\nabla\\Theta$ and~\\eqref{e30},~\\eqref{e14} that \\begin{align} &\\|\\hat{\\mathcal{D}}(k)-\\tilde{\\mathcal{D}}(k)\\|\\nonumber\\\\\n&\\leq \\sum_{i\\in V}(\\|\\hat{\\mathcal{D}}^{[i]}_x(k)-\\tilde{\\mathcal{D}}^{[i]}_x(k)\\| + \\|\\hat{\\mathcal{D}}^{[i]}_{\\mu}(k)-\\tilde{\\mathcal{D}}^{[i]}_{\\mu}(k)\\|)\\nonumber\\\\\n&\\leq 2N L_{\\Theta}\\sum_{\\tau=k-\\tau_{\\max}}^{k-1}\\|\\eta(\\tau+1)-\\eta(\\tau)\\|\\nonumber\\\\\n&\\leq 2N L_{\\Theta}\\sum_{\\tau=k-\\tau_{\\max}}^{k-1}\\alpha(\\tau)\\|\\mathcal{D}(\\tau)\\|.\\label{e45}\\end{align}\n\nSince $K$ is compact and the image of $\\nabla\\Theta$ is uniformly compact, ~\\eqref{e45} implies that there is $\\Upsilon > 0$ such that \\begin{align}&2\\alpha(k)\\|\\hat{\\mathcal{D}}(k)-\\tilde{\\mathcal{D}}(k)\\|\\|\\eta(k)-{\\eta}\\|\\nonumber\\\\\n&\\leq 2NL_{\\Theta}\\sum_{\\tau=k-\\tau_{\\max}}^{k-1}2\\alpha(k)\\alpha(\\tau)\n\\|\\mathcal{D}(\\tau)\\|\\|\\eta(k)-{\\eta}\\|\\nonumber\\\\\n&\\leq 2N L_{\\Theta}\\sum_{\\tau=k-\\tau_{\\max}}^{k-1}\\big(\\alpha(k)^2+\\alpha(\\tau)^2\n\\|\\mathcal{D}(\\tau)\\|^2\\|\\eta(k)-{\\eta}\\|^2\\big)\\nonumber\\\\\n&\\leq 2N L_{\\Theta}\\tau_{\\max}\\alpha(k)^2 + \\Upsilon\\sum_{\\tau=k-\\tau_{\\max}}^{k-1}\\alpha(\\tau)^2.\\label{e46}\\end{align}\n\nNow consider the term with $\\|\\tilde{\\mathcal{D}}(k)-\\mathcal{D}(k)\\|$ on the right-hand side in~\\eqref{e30}. Recall that $K$ is compact. By the Taylor expansion, we reach that \\begin{align}&\\mathcal{D}^{[i]}_x(k) =\\nabla_{x^{[i]}}f_i(x^{[i]}(k),\n\\Lambda_i(k))\\nonumber\\\\\n&+\\sum_{\\ell=1}^{m_i}\\mu^{[i]}_{\\ell}(k)\\nabla_{x^{[i]}}\nG^{[i]}_{\\ell}(x^{[i]}(k),\\Lambda_i(k))+O(c_i(k))\\textbf{1}_{n_i}.\\label{e21}\\end{align}\n\nWith the above relation, we have the following for $\\tilde{\\mathcal{D}}^{[i]}(k)-\\mathcal{D}^{[i]}(k)$:\n\\begin{align}\\tilde{\\mathcal{D}}^{[i]}_{\\mu}(k)-\\mathcal{D}^{[i]}_{\\mu}(k) = 0,\\quad\n\\tilde{\\mathcal{D}}^{[i]}_x(k)-\\mathcal{D}^{[i]}_x(k) = O(c_i(k))\\textbf{1}_{n_i}.\\label{e12}\\end{align}\n\nNotice that $\\|\\mathcal{D}(k)\\|$ is uniformly bounded. Substitute~\\eqref{e46} and~\\eqref{e12} into~\\eqref{e30}, sum over $[0,T]$, and it renders that there is some $\\Upsilon', \\Upsilon'' > 0$ such that the following estimate holds: \\begin{align}&2\\sum_{k=0}^T\\alpha(k)\\langle \\hat{\\mathcal{D}}(k), \\eta(k)-\\eta\\rangle\\nonumber\\\\\n&\\leq \\|\\eta(0) - \\eta\\|^2 - \\|\\eta(T+1) - \\eta\\|^2 + \\Upsilon'\\sum_{k=0}^T\\alpha(k)^2\\nonumber\\\\\n&+ \\Upsilon''\\sum_{k=0}^T\\alpha(k)\\max_{i\\in V}c_i(k).\\label{e26}\\end{align}\n\nSince $\\{\\alpha(k)^2\\}$ and $\\{\\alpha(k)\\max_{i\\in V}c_i(k)\\}$ are summable, then the right-hand side of~\\eqref{e26} is finite when we let $T\\rightarrow+\\infty$. We now show the following by contradiction: \\begin{align}\\liminf_{k\\rightarrow+\\infty}\\langle \\hat{\\mathcal{D}}(k),\\eta(k)-\\eta\\rangle\\leq0.\\label{e34}\\end{align}\nAssume that there are $k_0\\geq0$ and $\\epsilon>0$ such that $\\langle \\hat{\\mathcal{D}}(k),\\eta(k)-\\eta\\rangle \\geq \\epsilon$ for all $k\\geq k_0$. Since $\\alpha(k) > 0$ and $\\{\\alpha(k)\\}$ is not summable, then we have the following: \\begin{align}&\\sum_{k=0}^{+\\infty}\\alpha(k)\\langle \\hat{\\mathcal{D}}(k),\\eta(k)-\\eta\\rangle\\nonumber\\\\\n&\\geq \\sum_{k=0}^{k_0}\\alpha(k)\\langle \\hat{\\mathcal{D}}(k),\\eta(k)-\\eta\\rangle\\nonumber\\\\\n&+ \\sum_{k=k_0+1}^{+\\infty}\\alpha(k)\\langle \\hat{\\mathcal{D}}(k),\\eta(k)-\\eta\\rangle\\nonumber\\\\\n&\\geq \\sum_{k=0}^{k_0}\\alpha(k)\\langle \\hat{\\mathcal{D}}(k),\\eta(k)-\\eta\\rangle\n+ \\epsilon\\sum_{k=k_0+1}^{+\\infty}\\alpha(k)\\nonumber\\\\\n&\\geq +\\infty.\\nonumber\\end{align}\nWe then reach a contradiction, and thus~\\eqref{e34} holds. Equivalently, the following relation holds: \\begin{align*}\\limsup_{k\\rightarrow+\\infty}\\langle \\hat{\\mathcal{D}}(k),\\eta-\\eta(k)\\rangle\\geq0.\\end{align*} Since $\\nabla\\Theta$ is closed, the above relation implies that there is a limit point $\\tilde{\\eta}\\in K$ of the sequence $\\{\\eta(k)\\}$ such that the following holds: $\\langle \\nabla\\Theta(\\tilde{\\eta}),\\eta-\\tilde{\\eta}\\rangle\\geq0, \\quad \\forall\\eta\\in K$. Since $\\nabla\\Theta$ is pseudo-monotone, then we have the following relation: \\begin{align}\\langle \\nabla\\Theta(\\eta),\\eta-\\tilde{\\eta}\\rangle\\geq0, \\quad \\forall\\eta\\in K.\\label{e22}\\end{align}\n\nWe now set out to show the following by contradiction: \\begin{align}\\langle \\nabla\\Theta(\\tilde{\\eta}),\\eta-\\tilde{\\eta}\\rangle\\geq0,\\quad \\forall \\eta\\in K.\\label{e32}\\end{align} Assume that there is $\\hat{\\eta}\\in K$ such that the following holds: \\begin{align}\\langle \\nabla\\Theta(\\tilde{\\eta}),\\hat{\\eta}-\\tilde{\\eta}\\rangle < 0.\\label{e33}\\end{align} Now choose $\\varepsilon\\in(0,1)$, and define $\\eta_{\\varepsilon} \\triangleq \\hat{\\eta} + \\varepsilon(\\tilde{\\eta}-\\hat{\\eta})$. Since $K$ is convex, then we have $\\eta_{\\varepsilon}\\in K$. The following holds: \\begin{align}\\langle \\nabla\\Theta(\\eta_{\\varepsilon}),(1-\\varepsilon)(\\hat{\\eta}-\\tilde{\\eta})\\rangle = \\langle \\nabla\\Theta(\\eta_{\\varepsilon}),\\eta_{\\varepsilon}-\\tilde{\\eta}\\rangle\\geq0,\\label{e28}\n\\end{align} where in the inequality we use~\\eqref{e22}. It follows from~\\eqref{e28} that the following relation holds for any $\\varepsilon\\in(0,1)$: \\begin{align}\\langle \\nabla\\Theta(\\eta_{\\varepsilon}),\\hat{\\eta}-\\tilde{\\eta}\\rangle\\geq0,\\quad \\forall \\eta\\in K.\\label{e29}\n\\end{align} Since $\\nabla\\Theta$ is closed, letting $\\varepsilon\\rightarrow1$ in~\\eqref{e29} gives that: $\\langle \\nabla\\Theta(\\tilde{\\eta}),\\hat{\\eta}-\\tilde{\\eta}\\rangle\\geq0$, which contradicts~\\eqref{e33}. As a result, the relation~\\eqref{e32} holds. By~\\cite{Palomar.Eldar:10}, the relation~\\eqref{e32} implies that $\\tilde{\\eta}\\in\\mathbb{X}_{\\rm UC}$. We replay $\\eta$ by $\\eta(k)$ in~\\eqref{e32}, and establish the following:\n\\begin{align}\\langle \\nabla\\Theta(\\tilde{\\eta}),\\eta(k)-\\tilde{\\eta}\\rangle\\geq0.\\label{e31}\\end{align} Since $\\nabla\\Theta$ is pseudo-monotone, it follows from~\\eqref{e31} that\\begin{align}\\langle \\hat{\\mathcal{D}}(k), \\eta(k)-\\tilde{\\eta}\\rangle \\geq 0.\\label{e27}\\end{align} Replace $\\eta$ with $\\tilde{\\eta}$ in~\\eqref{e23}, apply~\\eqref{e27}, and it renders: \\begin{align}\\|\\eta(k+1) - \\tilde{\\eta}\\|^2 -\\|\\eta(k) - \\tilde{\\eta}\\|^2\n\\leq\\alpha(k)^2\\|\\mathcal{D}(k)\\|^2.\\label{e35}\\end{align} Sum up~\\eqref{e35} over $[s,k]$ and we have \\begin{align}\\|\\eta(k) - \\tilde{\\eta}\\|^2\n\\leq \\|\\eta(s) - \\tilde{\\eta}\\|^2 + \\sum_{\\tau=s}^{k-1}\\alpha(\\tau)^2\\|\\mathcal{D}(\\tau)\\|^2.\\label{e36}\\end{align}\nWe take the limits on $k$ first and then $s$ on both sides of~\\eqref{e36}. Since $\\{\\alpha(k)\\}$ is square summable and $\\{\\mathcal{D}(k)\\}$ is uniformly bounded, we have \\begin{align}\\limsup_{k\\rightarrow\\infty}\\|\\eta(k) - \\tilde{\\eta}\\|^2\n\\leq \\liminf_{s\\rightarrow\\infty}\\|\\eta(s) - \\tilde{\\eta}\\|^2.\\nonumber\\end{align} It implies that $\\{\\|\\eta(k) - \\tilde{\\eta}\\|\\}$ converges. Since $\\tilde{\\eta}$ is a limit point of $\\{\\eta(k)\\}$, $\\{\\eta(k)\\}$ converges to $\\tilde{\\eta}\\in\\mathbb{X}_{\\rm UC}$. Furthermore, we have $\\tilde{x}\\in\\mathbb{X}_{\\rm CVX}$ by (P2) in Proposition~\\ref{pro5}.\\oprocend\n\n\\subsection{Scenario two: strongly monotone reduced game map and known $\\tau_{\\max}$}\n\nIn the last section, the convergence of Algorithm~\\ref{ta:algo1} is ensured under the mild assumptions: $\\tau_{\\max}$ is unknown and the game map $\\nabla\\Theta$ is merely pseudo-monotone. This set of assumptions requires the usage of diminishing step-sizes. Note that diminishing step-sizes may cause a slow convergence rate. The shortcoming can be partially addressed by choosing a constant step-size when the inequality constraints are absent and (A5) in Theorem~\\ref{the1} is strengthened to the following one: \\begin{assumption} The map $\\nabla\\Theta^r$ is strongly monotone over $X$ with constant $\\rho$.\\label{asm5}\\end{assumption}\n\n\\begin{algorithm}[htbp]\\caption{The distributed gradient-based algorithm for Scenario two} \\label{ta:algo2}\n\n\\textbf{Require:} Each player in $V$ chooses the initial state $x^{[i]}(0)\\in X_i$.\n\n\\textbf{Ensure:} At each $k \\ge 0$, each player~$i\\in V$ executes the following steps:\n\n1. Each player $i\\in V$ updates its state according to the following rule:\\begin{align*} x^{[i]}(k+1) = \\mathbb{P}_{X_i}[x^{[i]}(k) - \\alpha D^{[i]}(k)],\\end{align*} where the approximate gradient $D^{[i]}(k)$ is given by: \\begin{align*}[D^{[i]}(k)]_{\\ell} &= \\frac{1}{2c_i(k)}\n\\{f_i(x^{[i]}(k)+c_i(k)e^{[\\ell]}_{n_i},\\Lambda_i(k))\\nonumber\\\\\n&-f_i(x^{[i]}(k)-c_i(k)e^{[\\ell]}_{n_i},\\Lambda_i(k))\\},\\end{align*} for $\\ell = 1,\\cdots,n_i$.\n\n2. Repeat for $k = k+1$.\n\\end{algorithm}\n\nAlgorithm~\\ref{ta:algo2} is proposed to address Scenario two. In particular, Algorithm~\\ref{ta:algo2} is similar to Algorithm~\\ref{ta:algo1}, but has two distinctions. Firstly, Algorithm~\\ref{ta:algo2} excludes the inequality constraints. Because the game map $\\nabla \\Theta$ cannot be strongly monotone due to $\\mathcal{H}_i$ being linear in $\\mu^{[i]}$. Secondly, thanks to the strong monotonicity of $\\nabla \\Theta$, a constant step-size replaces diminishing step-sizes. The following theorem summarizes the convergence properties of Algorithm~\\ref{ta:algo2}.\n\n\\begin{theorem} Suppose the following holds:\n\\begin{enumerate}\n\\item[(B1)] the quantity $0\\leq\\tau_{\\max}<\\frac{\\rho}{4NL_{\\Theta^r}}$ is known;\n\\item[(B2)] Assumptions~\\ref{asm7} and~\\ref{asm1} hold;\n\\item[(B3)] the step-size $\\alpha\\in(0,\\frac{\\rho-4NL_{\\Theta^r}\\tau_{\\max}}{L^2_{\\Theta^r}(2+16N^2\\tau_{\\max})})$;\n\\item[(B4)] The sequence $\\{\\max_{i\\in V}c_i(k)\\}$ is summable.\n\\item[(B5)] Assumption~\\ref{asm5} holds.\n\\end{enumerate}\nFor any initial state $x(0)\\in X$, the sequence of $\\{x(k)\\}$ generated by Algorithm~\\ref{ta:algo2} converges to $\\tilde{x}\\in\\mathbb{X}_{\\rm CVX}$.\\label{the2}\n\\end{theorem}\n\n\\textbf{Proof:} First of all, for any $k\\geq0$ and $i$, we define $\\hat{\\mathcal{D}}^{[i]}(k) \\triangleq \\nabla_{x^{[i]}}f_i(x^{[i]}(k),\\{x^{[j]}(k)\\}_{j\\in\\mathcal{N}^{\\rm IN}_i})$, which is the gradient evaluated at the delay-free states. So, the quantity $\\hat{\\mathcal{D}}(k)$ is free of time delays and perturbations. Similarly, for any $k\\geq0$ and $i$, we define $\\tilde{\\mathcal{D}}^{[i]}(k) \\triangleq \\nabla_{x^{[i]}}f_i(x^{[i]}(k),\\Lambda_i(k))$ which is the gradient evaluated at the delayed states. Then the quantity $\\tilde{\\mathcal{D}}(k)$ is free of perturbations but subject to time delays.\n\nWe then write Algorithm~\\ref{ta:algo2} in the following compact form: \\begin{align}x(k+1) = \\mathbb{P}_X[x(k)-\\alpha D(k)].\\label{e39}\n\\end{align} It is noticed that $\\tilde{x}\\in{\\mathbb{X}}_{\\rm CVX}$ is a fixed point of the operator $\\mathbb{P}_X[\\cdot-\\alpha \\nabla\\Theta^r(\\cdot)]$ for any $\\alpha > 0$; i.e., it holds that $\\tilde{x} = \\mathbb{P}_X[\\tilde{x}-\\alpha \\nabla\\Theta^r(\\tilde{x})]$. By this, we have the following relations: \\begin{align}&\\|x(k+1)-\\tilde{x}\\|^2\\nonumber\\\\\n&= \\|\\mathbb{P}_X[x(k)-\\alpha D(k)] - \\mathbb{P}_X[\\tilde{x}-\\alpha \\nabla\\Theta^r(\\tilde{x})]\\|^2\\nonumber\\\\\n&\\leq \\|(x(k) - \\tilde{x})-\\alpha(D(k)-\\nabla\\Theta^r(\\tilde{x}))\\|^2\\nonumber\\\\\n&= \\|x(k) - \\tilde{x}\\|^2-2\\alpha\\langle x(k) - \\tilde{x}, \\hat{D}(k)-\\nabla\\Theta^r(\\tilde{x})\\rangle\\nonumber\\\\\n&+ \\alpha^2\\|(D(k)-\\tilde{D}(k))+(\\tilde{D}(k)-\\hat{D}(k))+(\\hat{D}(k)-u)\\|^2\\nonumber\\\\\n&-2\\alpha\\langle x(k) - \\tilde{x}, D(k)-\\tilde{D}(k)\\rangle\\nonumber\\\\\n&-2\\alpha\\langle x(k) - \\tilde{x}, \\tilde{D}(k)-\\hat{D}(k)\\rangle\\nonumber\\\\\n&\\leq (1-2\\rho\\alpha+4L_{\\Theta^r}^2\\alpha^2)\\|x(k)-\\tilde{x}\\|^2\\nonumber\\\\\n&+ 4\\alpha^2(\\|D(k)-\\tilde{D}(k)\\|^2+\\|\\tilde{D}(k)-\\hat{D}(k)\\|^2)\\nonumber\\\\\n&+2\\alpha\\|x(k) - \\tilde{x}\\|(\\|D(k)-\\tilde{D}(k)\\|+\\|\\tilde{D}(k)-\\hat{D}(k)\\|),\\label{e42}\\end{align}\nwhere we use the non-expansiveness property of the projection operator $\\mathbb{P}_X$ in the first inequality, and the strong monotonicity and the Lipschitz continuity of $\\nabla\\Theta^r$ in the last inequality. For the term of $\\|\\hat{D}(k)-\\tilde{D}(k)\\|$ in~\\eqref{e42}, one can derive the following relations: \\begin{align}&\\|\\hat{D}(k)-\\tilde{D}(k)\\| \\leq NL_{\\Theta^r}\\sum_{\\tau=k-\\tau_{\\max}}^{k-1}\\|x(\\tau+1)-x(\\tau)\\|\\nonumber\\\\\n&\\leq NL_{\\Theta^r}\\sum_{\\tau=k-\\tau_{\\max}}^{k-1}\\big(\\|x(\\tau+1)-\\tilde{x}\\|\n+\\|x(\\tau)-\\tilde{x}\\|\\big)\\nonumber\\\\\n&\\leq 2NL_{\\Theta^r}\\sum_{\\tau=k-\\tau_{\\max}}^{k-1}\\|x(\\tau)-\\tilde{x}\\|.\\label{e47}\\end{align}\n\nSimilar to~\\eqref{e12}, we have \\begin{align}\\tilde{D}^{[i]}(k)-D^{[i]}(k) = O(c_i(k))\\textbf{1}_{n_i}.\\label{e13}\\end{align}\n\nSubstituting~\\eqref{e47} and~\\eqref{e13} into~\\eqref{e42} yields: \\begin{align}&\\|x(k+1)-\\tilde{x}\\|^2\\nonumber\\\\\n&\\leq (1-2\\rho\\alpha+4L_{\\Theta^r}^2\\alpha^2+4NL_{\\Theta^r}\\alpha\\tau_{\\max})\n\\|x(k)-\\tilde{x}\\|^2\\nonumber\\\\\n&+ (32N^2L_{\\Theta^r}^2\\alpha^2+4NL_{\\Theta^r}\\alpha)\n\\sum_{\\tau=k-\\tau_{\\max}}^{k-1}\\|x(\\tau)-\\tilde{x}\\|^2\\nonumber\\\\\n&+4\\alpha^2(\\max_{i\\in V}c_i(k))^2 + 2\\alpha\\max_{i\\in V}c_i(k)\\|x(k)-\\tilde{x}\\|,\\label{e48}\\end{align}\nwhere we use the relation $2ab \\leq a^2 + b^2$.\n\nTo study the convergence of $x(k)$ in~\\eqref{e48}, we define the following notation: \\begin{align*}&\\psi(k) \\triangleq x(k) - \\tilde{x},\\quad\\|\\chi(k)\\| \\triangleq \\max_{k-\\tau_{\\max}\\leq \\tau \\leq k-1}\\|\\psi(\\tau)\\|,\\nonumber\\\\\n&e(k)\\triangleq 4\\alpha^2(\\max_{i\\in V}c_i(k))^2 + 2\\alpha\\max_{i\\in V}c_i(k)\\|x(k)-\\tilde{x}\\|.\\end{align*} Then~\\eqref{e48} is rewritten as follows: \\begin{align}\\|\\psi(k+1)\\|\\leq a\\|\\psi(k)\\|+b\\|\\chi(k)\\|+e(k).\\label{e43}\\end{align} where $a \\triangleq (1-2\\rho\\alpha+4L_{\\Theta^r}^2\\alpha^2+4NL_{\\Theta^r}\\alpha\\tau_{\\max})$, $b \\triangleq \\tau_{\\max}(32N^2L_{\\Theta^r}^2\\alpha^2+4NL_{\\Theta^r}\\alpha)$ and the sequence $\\{e(k)\\}$ is diminishing. From the recursion of~\\eqref{e43}, we can derive the following relation for any pair of $k>s\\geq0$: \\begin{align}&\\|\\psi(k)\\| \\leq a^{k-s}\\|\\psi(s)\\| + \\sum_{\\tau=s}^{k-1}a^{k-\\tau}b\\chi(\\tau) + \\sum_{\\tau=s}^{k-1}e(\\tau)\\nonumber\\\\\n&\\leq a^{k-s}\\|\\psi(s)\\| + b\\sup_{s\\leq\\tau\\leq k}\\|\\chi(\\tau)\\|\\sum_{\\tau=s}^{k-1}a^{k-\\tau} + \\sum_{\\tau=s}^{k-1}e(\\tau)\\nonumber\\\\\n&\\leq a^{k-s}\\|\\psi(s)\\| + \\frac{b}{1-a}\\sup_{s\\leq\\tau\\leq k}\\|\\chi(\\tau)\\| + \\sum_{\\tau=s}^{k-1}e(\\tau).\\label{e44}\\end{align} Recall that $\\{\\max_{i\\in V}c_i(k)\\}$ and thus $\\{(\\max_{i\\in V}c_i(k))^2\\}$ are summable. Take the limits on $k$ and $s$ at both sides of~\\eqref{e44}, and it gives the following relation: \\begin{align}\\limsup_{k\\rightarrow+\\infty}\\|\\psi(k)\\| \\leq \\frac{b}{1-a}\\limsup_{k\\rightarrow+\\infty}\\|\\chi(k)\\|.\\label{e49}\\end{align} On the other hand, one can see that \\begin{align}\\limsup_{k\\rightarrow+\\infty}\\|\\chi(k)\\| = \\limsup_{k\\rightarrow+\\infty}\\|\\psi(k)\\|.\\label{e50}\\end{align}\n\nSince $\\alpha\\in(0,\\frac{\\rho-4NL_{\\Theta^r}\\tau_{\\max}}{L^2_{\\Theta^r}(2+16N^2\\tau_{\\max})})$, then $\\frac{a}{1-b}<1$. The combination of~\\eqref{e49} and~\\eqref{e50} renders that \\begin{align}\\limsup_{k\\rightarrow+\\infty}\\|\\psi(k)\\| = 0.\\label{e40}\\end{align} Apparently, \\begin{align}\\liminf_{k\\rightarrow+\\infty}\\|\\psi(k)\\| \\geq 0.\\label{e41}\\end{align} The combination of~\\eqref{e40} and~\\eqref{e41} establishes the convergence of $\\{x(k)\\}$ to $\\tilde{x}$. It completes the proof.\\oprocend\n\n\\section{Discussions}\\label{sec:discussion}\n\n\\subsection{Comparison of two scenarios}\\label{sec:comparison}\n\nThe two proposed algorithms are complementary. Algorithm~\\ref{ta:algo1} can address inequality constraints, does not need to know $\\tau_{\\max}$ and merely requires the game map to be pseudo monotone. It comes with the price of potentially slow convergence due to the utilization of diminishing step-sizes. In contrast, Algorithm~\\ref{ta:algo2} cannot deal with inequality constraints, needs to know $\\tau_{\\max}$ and requires the game map to be strongly monotone. It comes with the benefit of potentially fast convergence due to the utilization of a constant step-size.\n\n\n\\subsection{Discussion on Assumption~\\ref{asm1}}\\label{sec:assumption}\n\nThe proposed algorithms rely upon Assumption~\\ref{asm1}. In what follows, we will provide two sufficient conditions for Assumption~\\ref{asm1}. Recall that $\\sigma_{i,\\min}$, $\\sigma_{i,\\max}$,\n$\\sigma_{\\min}$ and $\\sigma_{\\max}$ are defined in Section~\\ref{subsection:notations}.\n\n\\subsubsection{Global Slater vectors}\n\n\\begin{assumption} There exist $\\bar{x}\\in X$ and $\\sigma' > 0$ such that $\\|-G^{[i]}(\\bar{x})\\|_{\\infty} \\geq \\sigma'$ for all $i\\in V$.\\label{asm3}\n\\end{assumption}\n\nThe vector $\\bar{x}$ that satisfies Assumption~\\ref{asm3} is referred to as the \\emph{global} Slater vector. The existence of the global Slater vector ensures the boundedness of $\\mathbb{X}_{\\rm UC}(M)$ as follows.\n\n\\begin{lemma} If Assumption~\\ref{asm3} holds and $G^{[i]}_{\\ell}$ is convex in $x$, then for any $(\\tilde{x},\\tilde{\\mu})\\in \\mathbb{X}_{\\rm UC}(\\mathds{R}^m_{\\geq0})$, it holds that $\\|\\tilde{\\mu}^{[i]}\\|_{\\infty} \\leq \\frac{\\sigma_{\\max}-\\sigma_{\\min}}{\\sigma'}$ for $i\\in V$.\\label{lem2}\n\\end{lemma}\n\n\\textbf{Proof:} Pick any $\\tilde{\\eta}\\triangleq(\\tilde{x},\\tilde{\\mu})\\in \\mathbb{X}_{\\rm UC}(\\mathds{R}^m_{\\geq0})$. It holds that \\begin{align}\\langle \\nabla\\Theta(\\tilde{\\eta}),\\tilde{\\eta}-\\eta\\rangle\\leq0.\\label{e55}\\end{align}\nChoose $\\eta = (\\bar{x}^T,0^T)^T$ in~\\eqref{e55}, and we have\n\\begin{align}&0\\geq\\langle \\nabla\\Theta^r(\\tilde{x}),\\tilde{x}-\\bar{x}\\rangle - \\sum_{i\\in V}\\langle \\tilde{\\mu}^{[i]},G^{[i]}(\\tilde{x})\\rangle\\nonumber\\\\\n&+ \\sum_{i\\in V}\\sum_{\\ell=1}^{m_i}\\tilde{\\mu}_{\\ell}^{[i]}\\langle \\nabla_{x^{[i]}} G^{[i]}_{\\ell}(\\tilde{x}),\\tilde{x}^{[i]}-\\bar{x}^{[i]}\\rangle\\nonumber\\\\\n&\\geq \\langle \\nabla\\Theta^r(\\tilde{x}),\\tilde{x}-\\bar{x}\\rangle - \\sum_{i\\in V}\\langle \\tilde{\\mu}^{[i]},G^{[i]}(\\bar{x})\\rangle,\\label{e56}\n\\end{align} where the last inequality uses $\\tilde{\\mu}^{[i]}_{\\ell}\\geq0$ and Taylor theorem and the convexity of $G^{[i]}_{\\ell}$ in $x$. It follows from~\\eqref{e56} and Assumption~\\ref{asm3} that \\begin{align}\\|\\tilde{\\mu}^{[i]}\\|_{\\infty}\\leq\\frac{-\\langle \\nabla\\Theta^r(\\tilde{x}),\\tilde{x}-\\bar{x}\\rangle}{\\sigma'}\\leq \\frac{\\sigma_{\\max}-\\sigma_{\\min}}{\\sigma'}.\\nonumber\\end{align}\n\\oprocend\n\n\n\\subsubsection{Private Slater vectors}\n\nIn Lemma~\\ref{lem2}, each player has to access the global information of $\\bar{x}$, $\\sigma'$, $\\sigma_{\\min}$ and $\\sigma_{\\max}$. In what follows, we will derive a sufficient condition which only requires private information of each player~$i$ to determine an upper bound on $\\mu^{[i]}$.\n\n\\begin{assumption} For each $i\\in V$, there exists $\\sigma_i > 0$ and $\\bar{x}_i\\in X_i$ such that for any $(\\tilde{x},\\tilde{\\mu})\\in\\mathbb{X}_{\\rm UC}$, $\\|-G^{[i]}(\\bar{x}^{[i]},\\tilde{x}^{[-i]})\\|_{\\infty} \\geq \\sigma_i$ holds.\\label{asm2}\n\\end{assumption}\n\nThe vector $\\bar{x}^{[i]}$ that satisfies Assumption~\\ref{asm3} is referred to as the \\emph{private} Slater vector. One case where Assumption~\\ref{asm2} holds is that $G^{[i]}$ only depends upon $x^{[i]}$ and, for each $i\\in V$, there is $\\bar{x}^{[i]}\\in X_i$ such that $\\|-G^{[i]}(\\bar{x}^{[i]})\\|_{\\infty} \\geq \\sigma_i$ holds for some $\\sigma_i > 0$.\n\n\\begin{lemma} If Assumption~\\ref{asm2} holds, then for any $(\\tilde{x},\\tilde{\\mu})\\in \\mathbb{X}_{\\rm UC}(\\mathds{R}^m_{\\geq0})$, it holds that $\\|\\tilde{\\mu}^{[i]}\\|_{\\infty} \\leq \\frac{\\sigma_{i,\\max}-\\sigma_{i,\\min}}{\\sigma_i}$ for $i\\in V$.\\label{lem5}\n\\end{lemma}\n\n\\textbf{Proof:} Pick any $(\\tilde{x},\\tilde{\\mu})\\in\\mathbb{X}_{\\rm UC}(\\mathds{R}^m_{\\geq0})$. By Assumption~\\ref{asm2}, there is $\\bar{x}^{[i]}\\in X_i$ such that \\begin{align*}\\|-G^{[i]}(\\bar{x}^{[i]},\\tilde{x}^{[-i]})\\|_{\\infty} \\geq \\sigma.\\end{align*} Since $(\\tilde{x},\\tilde{\\mu})\\in\\mathbb{X}_{\\rm UC}$, we then have the following: \\begin{align}&\\mathcal{H}_i(\\tilde{x}^{[i]},\\tilde{x}^{[-i]},0)\\leq\n\\mathcal{H}_i(\\tilde{x}^{[i]},\\tilde{x}^{[-i]},\\tilde{\\mu}^{[i]})\\nonumber\\\\\n&\\leq\n\\mathcal{H}_i(\\bar{x}^{[i]},\\tilde{x}^{[-i]},\\tilde{\\mu}^{[i]})\\nonumber\\end{align} where in the first inequality we use $0\\in M_i$ and in the second inequality we use $\\bar{x}^{[i]}\\in X_i$. The above relations further render the following: \\begin{align}f_i(\\tilde{x})\\leq f_i(\\bar{x}^{[i]},\\tilde{x}^{[-i]}) + \\langle \\tilde{\\mu}^{[i]}, G^{[i]}(\\bar{x}^{[i]},\\tilde{x}^{[-i]})\\rangle.\\label{e24}\\end{align} Recall $\\|-G^{[i]}(\\bar{x}^{[i]},\\tilde{x}^{[-i]})\\|_{\\infty} \\geq \\sigma_i$. Then the relation~\\eqref{e24} implies the relation of $\\|\\tilde{\\mu}^{[i]}\\|_{\\infty} \\leq \\frac{\\sigma_{i,\\max}-\\sigma_{i,\\min}}{\\sigma_i}$.\\oprocend\n\n\n\n\n\n\n\\subsection{Future directions}\n\nThe current paper imposes several assumptions: (1) the compactness of $X_i$; (2) the smoothness of component functions; (3) identical $\\alpha(k)$ or $\\alpha$ for all the players; (4) the convexity of component functions and constraints. In particular, the compactness of $X_i$ ensures the uniform boundedness of partial gradients. The convexity guarantees that local optimum are also globally optimal in the sense of Nash equilibrium. It is of interest to relax these assumptions.\n\n\\section{Numerical simulations}\\label{sec:simulations}\n\nIn this section, we will provide a set of numerical simulations to verify the performance of our proposed algorithms. The procedure ${\\rm sample}(A)$ returns a uniform sample from the set $A$.\n\n\\subsection{Algorithm~\\ref{ta:algo1}}\\label{sec:simulations-alg1}\n\nConsider a power network which can be modeled as an interconnected graph $\\mathcal{G}_p\\triangleq \\{\\mathcal{V}_p,\\mathcal{E}_p\\}$ where each node $i\\in \\mathcal{V}_p$ represents a bus and each link in $\\mathcal{E}_p$ represents a branch. The buses in $\\mathcal{V}_p$ are indexed by $1,\\cdots,N$ and bus $1$ denotes the feeder which has a fixed voltage and flexible power injection. A lossless DC model is used to characterize the relation between power generations and loads at different buses and power flows across various branches. Each bus is connected to either a power load or supply and each load is associated with an end-user. The set of load buses is denoted by $\\mathcal{V}_l\\subseteq \\mathcal{V}_p$ and the set of supply buses is denoted by $\\mathcal{V}_s\\subseteq \\mathcal{V}_p$.\n\nThe maximum available power supply at bus $i\\in \\mathcal{V}_s$ is denoted by $S_i\\geq0$, and the intended power load at bus $i\\in \\mathcal{V}_l$ is denoted by $L_i\\geq0$. If the total power supply exceeds the total intended power load, then all the intended loads can be satisfied. Otherwise, some loads need to reduce. The reduced load of end-user~$i$ is denoted by $R_i\\in[0,L_i]$. The objective function of each end-user~$i$ is given by $f_i(R_i) = c_i R_i + p(\\textbf{1}^T_{|\\mathcal{V}_l|}(L-R))(L_i-R_i) - u_i(L_i-R_i)$. The quantity $c_iR_i$ represents the disutility induced by load reduction $R_i$ with $c_i>0$. The scalar $p(\\textbf{1}^T(L-R))$ is the charged price given the total actual load $\\textbf{1}^T(L-R)$. The value $u_i(L_i-R_i)$ stands for the benefit produced by load $L_i-R_i$.\n\nEach end-user~$i$ aims to minimize its own objective function $f_i(R_i)$ by adjusting $R_i$. Such decision making is subject to the physical constraints of the power grid. The first constraint is that the total actual load cannot exceed the maximum available supply; i.e., \\begin{align}\n\\textbf{1}_{|\\mathcal{V}_l|}^T(L-R) \\leq \\textbf{1}_{|\\mathcal{V}_s|}^TS.\\label{e52}\\end{align} Another set of constraints are induced by the limitations of power flows on branches. Let $H_g\\in[-1,1]^{|\\mathcal{E}|\\times|\\mathcal{V}_g|}$ (resp. $H_l\\in[-1,1]^{|\\mathcal{E}|\\times|\\mathcal{V}_l|}$) be the generation (resp. load) shift factor matrix. For $H_g$, the $(i,\\ell)$ entry represents the power that is distributed on line $\\ell$ when $1MW$ is injected into bus $i$ and withdrawn at the reference bus. Denote by $f_e^{\\max}$ the maximum capacity of branch $e$. The power flow constraint for branch $e$ can be expressed as: \\begin{align}-f^{\\max}\\leq H_g S - H_l (L-R)\\leq f^{\\max},\\label{e53}\\end{align} where $f^{\\max}\\triangleq{\\rm col}[f_e^{\\max}]_{e\\in \\mathcal{E}_p}$, $S\\triangleq {\\rm col}[S_i]_{i\\in\\mathcal{V}_s}$, $L\\triangleq {\\rm col}[L_i]_{i\\in\\mathcal{V}_l}$ and $R\\triangleq {\\rm col}[R_i]_{i\\in\\mathcal{V}_l}$.\n\nThe above description defines a non-cooperative game among end-users in $\\mathcal{V}_l$ as follows:\n\\begin{align}&\\min_{R_i\\in [0,L_i]}c_i R_i + p(\\textbf{1}_{|\\mathcal{V}_l|}^T(L-R))(L_i-R_i) - u_i(L_i-R_i)\\nonumber\\\\\n&{\\rm s.t.}\\quad \\textbf{1}_{|\\mathcal{V}_l|}^T(L-R) \\leq \\textbf{1}_{|\\mathcal{V}_s|}^TS\\nonumber\\\\\n&\\quad\\quad H_g S - H_l (L-R) - f^{\\max}\\leq 0\\nonumber\\\\\n&\\quad\\quad - H_g S + H_l (L-R) -f^{\\max}\\leq 0.\\label{e54}\n\\end{align}\n\nIn game~\\eqref{e54}, $p$ is the pricing policy of load serving entity (LSE). This is confidential information of LSE and should not be disclosed to end-users. So the objective function $f_i$ is partially unknown to end-user~$i$. In addition, the numerical values of generation and load shift factor matrices are of national security interest and are kept confidential from the public. Therefore, the power flow constraints are unknown to end-users. Assume that all the end-users can communicate with LSE. Given $R$, LSE broadcasts the price value $p(\\textbf{1}_{|\\mathcal{V}_l|}^T(L-R))$, the power imbalance $\\textbf{1}_{|\\mathcal{V}_l|}^T(L-R) - \\textbf{1}_{|\\mathcal{V}_s|}^TS$, and power limit violations $H_g S - H_l (L-R) - f^{\\max}$ and $- H_g S + H_l (L-R) -f^{\\max}$. In this way, end-users can access the values of $f_i$ and $G^{[i]}$ in Algorithm~\\ref{ta:algo1} without knowing their structures.\\footnote{Many networked engineering systems operate in a hierarchical structure; e.g., Internet, power grid and transportation systems, where a central system operator is placed at the top layer and end-users are placed at the bottom layer~\\cite{MC-SHL-ARC-JCD:07a,Wu.Moslehi.Bose:89}. Here we assume that LSE can communicate with all the end-users. This assumption is widely used in; e.g., network flow control~\\cite{LL-99}.}\n\nWe choose $u_i(L_i-R_i) = -\\frac{1}{2}a_i(L_i-R_i)^2+b_i(L_i-R_i)$ with $a_i>0$. For the pricing mechanism, we use (4) in~\\cite{Bulow.Peiderer:83}; i.e., $p(w) = \\beta w^{\\frac{1}{\\eta}}$ with $\\beta>0$, and $\\eta\\in(0,1)$. So $f_i(R) = c_iR_i + (L_i-R_i)\\beta (\\textbf{1}_{|\\mathcal{V}_l|}^T(L-R))^{\\frac{1}{\\eta}} +\\frac{1}{2}a_i(L_i-R_i)^2-b_i(L_i-R_i)$. One can compute the following: \\begin{align}\\frac{\\partial^2f_i}{\\partial^2R_i}&=a_i-2\\beta\\frac{1}{\\eta}\n(\\textbf{1}_{|\\mathcal{V}_l|}^T(L-R))^{\\frac{1}{\\eta}-1}\\nonumber\\\\\n&-\\beta (L_i-R_i)\\frac{1}{\\eta}(\\frac{1}{\\eta}-1)(\\textbf{1}_{|\\mathcal{V}_l|}^T(L-R))^{\\frac{1}{\\eta}-2},\\nonumber\\\\\n\\frac{\\partial^2f_i}{\\partial R_i \\partial R_j}&=-\\beta\\frac{1}{\\eta}\n(\\textbf{1}_{|\\mathcal{V}_l|}^T(L-R))^{\\frac{1}{\\eta}-1}\\nonumber\\\\\n&-\\beta (L_i-R_i)\\frac{1}{\\eta}(\\frac{1}{\\eta}-1)(\\textbf{1}_{|\\mathcal{V}_l|}^T(L-R))^{\\frac{1}{\\eta}-2}.\\nonumber\\end{align}\n\nChoose $a_i$ such that $a_i > (N-2)\\beta\\frac{1}{\\eta}\n(\\textbf{1}_{|\\mathcal{V}_l|}^TL)^{\\frac{1}{\\eta}-1}+(N-1)\\beta L_i\\frac{1}{\\eta}(\\frac{1}{\\eta}-1)(\\textbf{1}_{|\\mathcal{V}_l|}^TL)^{\\frac{1}{\\eta}-2}>0$. Then the monotonicity property holds. Assume that there is a global Slater vector. Figure~\\ref{fig1} shows the simulation results for the IEEE $30$-bus Test System~\\cite{PSTCA} and the system parameters are adopted from MATPOWER~\\cite{MATPOWER}. The delays at each iteration are randomly chosen from $0$ to $10$.\n\n\\begin{figure*}[bh]\n \\centering\n \\includegraphics[width = \\linewidth]{DemandResponse}\n \\caption{The estimates of Algorithm~\\ref{ta:algo1} at $30$ buses for $\\tau_{\\max} = 10$.}\\label{fig1}\n\\end{figure*}\n\n\n\\subsection{Algorithm~\\ref{ta:algo2}}\\label{sec:simulations-alg2}\n\nConsider $10$ players and their components functions are defined as follows: \\begin{align*}f_i(x) = \\frac{1}{2}((Q^{[i]})^Tx)^2 + (q^{[i]})^Tx,\\end{align*} where\n\\begin{itemize}\n\\item $X_i = [\\psi_i^m\\;\\;\\psi_i^M]$ with $\\psi_i^m = {\\rm sample}([-10\\;\\;-5])$ and $\\psi_i^M = {\\rm sample}([5\\;\\;10])$;\n\\item $Q^{[i]}\\in\\mathds{R}^{10}$ with $Q^{[i]}_i = 1$ and $Q^{[i]}_j = \\frac{1}{10}$ for $j\\neq i$;\n\\item $q^{[i]}\\in\\mathds{R}^{10}$ with $q^{[i]}_j = {\\rm sample}([-50\\;\\;50])$ for $j\\in V$.\n\\end{itemize}\n\nOne can verify that $\\rho = \\frac{1}{10}$ and $L_{\\Theta^r} = 2$ in Theorem~\\ref{the2}. By (B4), we have $\\varrho_b>1$. In the simulation, we choose $\\varrho_b = 3$. Figure~\\ref{fig4} presents the estimate evolution of the players.\n\n\\begin{figure}[h]\n \\centering\n \\includegraphics[width = \\linewidth]{Algorithm2-x}\n \\caption{The estimates of Algorithm~\\ref{ta:algo2} for $\\tau_{\\max} = 10$.}\\label{fig4}\n\\end{figure}\n\n\\section{Conclusions}\\label{sec:conclusions}\n\nWe have studied a set of distributed robust adaptive computation algorithms for a class of generalized convex games. We have shown their asymptotic convergence to Nash equilibrium in the presence of network unreliability and the uncertainties of component functions and illustrated the algorithm performance via numerical simulations.\n\n\n\n\\bibliographystyle{plain}\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzehft b/data_all_eng_slimpj/shuffled/split2/finalzzehft new file mode 100644 index 0000000000000000000000000000000000000000..49176b7e97943c68f40063f837b4b292f902561b --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzehft @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction} \\label{sec:intro}\n\n\\IEEEPARstart{L}{earning} convolutional operators from large datasets is a growing trend in signal\/image processing, computer vision, and machine learning.\nThe widely known \\emph{patch-domain} approaches for learning kernels (e.g., filter, dictionary, frame, and transform) extract patches from training signals for simple mathematical formulation and optimization, yielding (sparse) features of training signals \\cite{Aharon&Elad&Bruckstein:06TSP, Elad&Aharon:06TIP, Yaghoobi&etal:13TSP, Hawe&Kleinsteuber&Diepold:13TIP, Mairal&Bach&Ponce:14FTCGV, Cai&etal:14ACHA, Ravishankar&Bressler:15TSP, Pfister&Bresler:15SPIE, Coates&NG:bookCh}.\nDue to memory demands, using many overlapping patches across the training signals hinders using large datasets\nand building hierarchies on the features, \ne.g., deconvolutional neural networks \\cite{Zeiler&etal:10CVPR}, \nconvolutional neural network (CNN) \\cite{LeCun&etal:98ProcIEEE}, \nand multi-layer convolutional sparse coding \\cite{Papyan&Romano&Elad:17JMLR}.\nFor similar reasons, the memory requirement of patch-domain approaches discourages learned kernels from being applied to large-scale inverse problems. \n\n\n\n\nTo moderate these limitations of the patch-domain approach, the so-called \\emph{convolution} perspective has been recently introduced by learning filters and obtaining (sparse) representations directly from the original signals without storing many overlapping patches, e.g., convolutional dictionary learning (CDL) \\cite{Zeiler&etal:10CVPR, Bristow&etal:13CVPR, Heide&eta:15CVPR, Wohlberg:16TIP, Chun&Fessler:18TIP, Chun&Fessler:17SAMPTA}. \nFor large datasets, CDL using careful algorithmic designs \\cite{Chun&Fessler:18TIP} is more suitable for learning filters than patch-domain dictionary learning \\cite{Aharon&Elad&Bruckstein:06TSP}; in addition, CDL can learn translation-invariant filters without obtaining highly redundant sparse representations \\cite{Chun&Fessler:18TIP}.\nThe CDL method applies the convolution perspective for learning kernels within \\dquotes{synthesis} signal models.\nWithin \\dquotes{analysis} signal models, however, there exist no prior frameworks using the convolution perspective for learning convolutional operators, whereas patch-domain approaches for learning analysis kernels are introduced in \\cite{Yaghoobi&etal:13TSP, Hawe&Kleinsteuber&Diepold:13TIP, Cai&etal:14ACHA, Ravishankar&Bressler:15TSP, Pfister&Bresler:15SPIE}.\n(See brief descriptions about synthesis and analysis signal models in \\cite[Sec.~\\Romnum{1}]{Hawe&Kleinsteuber&Diepold:13TIP}.)\n\n\n\nResearchers interested in dictionary learning have actively studied the structures of kernels learned by the patch-domain approach \\cite{Yaghoobi&etal:13TSP, Hawe&Kleinsteuber&Diepold:13TIP, Cai&etal:14ACHA, Ravishankar&Bressler:15TSP, Pfister&Bresler:15SPIE, Barchiesi&Plumbley:13TSP, Bao&Cai&Ji:13ICCV, Ravishankar&Bressler:13ICASSP}. \nIn training CNNs (see Appendix~\\ref{sec:CNN}), however,\nthere has been less study of filter structures having non-convex constraints, e.g., orthogonality and unit-norm constraints in Section~\\ref{sec:CAOL}, although it is thought that diverse (i.e., incoherent) filters can improve performance for some applications, e.g., image recognition \\cite{Coates&NG:bookCh}.\nOn the application side, researchers have applied (deep) NNs to signal\/image recovery problems.\nRecent works combined model-based image reconstruction (MBIR) algorithm with image refining networks \\cite{Yang&etal:16NIPS, Zhang&etal:17CVPR, Chen&Pock:17PAMI, Chen&etal:17arXiv, Wu&etal:17Fully3D, Romano&Elad&Milanfar:17SJIS, Buzzard&etal:18SJIS, Chun&Fessler:18IVMSP, Chun&etal:18arXiv:momnet, Chun&etal:18Allerton}.\nIn these \\emph{iterative} NN methods, \nrefining NNs should satisfy the non-expansiveness for fixed-point convergence \n\\cite{Chun&etal:18arXiv:momnet}; \nhowever, their trainings lack consideration of filter diversity constraints, e.g., orthogonality constraint in Section~\\ref{sec:CAOL}, \nand thus it is unclear whether the trained NNs are nonexpansive mapping \\cite{Chun&etal:18Allerton}.\n\n\n\n\\begin{figure*}[!pt]\n\\centering\n\n\\begin{tabular}{c}\n\\includegraphics[trim={3cm 11.2cm 3cm 11.3cm},clip,scale=0.46]{.\/Fig_v1\/diagram_abstract.pdf}\n\\end{tabular}\n\n\\vspace{-0.25em}\n\\caption{A general flowchart from\nlearning sparsifying operators $\\cO$ to \nsolving inverse problems via MBIR using learned operators $\\cO^\\star$;\nsee Section~\\ref{sec:back}.\nFor the $l\\rth$ training sample $x_l$, \n$F(\\cO; x_l)$ measures its sparse representation or sparsification errors, and\nsparsity of its representation generated by $\\cO$.\n}\n\\label{diag:abstract}\n\\end{figure*}\n\n\n\n\nThis paper proposes \\textit{1)} a new \\textit{convolutional analysis operator learning} (CAOL) framework that learns an analysis sparsifying regularizer with the convolution perspective, and \\textit{2)} a new convergent \\textit{Block Proximal Extrapolated Gradient method using a Majorizer} (BPEG-M \\cite{Chun&Fessler:18TIP}) for solving block multi-nonconvex problems \\cite{Xu&Yin:17JSC}.\nTo learn diverse filters, we propose \\textit{a)} CAOL with an orthogonality constraint that enforces a tight-frame (TF) filter condition in convolutional perspectives, and \\textit{b)} CAOL with a regularizer that promotes filter diversity.\nBPEG-M with sharper majorizers converges significantly faster than the state-of-the-art technique, Block Proximal Gradient (BPG) method \\cite{Xu&Yin:17JSC} for CAOL.\nThis paper also introduces a new X-ray computational tomography (CT) MBIR model using a convolutional sparsifying regularizer learned via CAOL \\cite{Chun&Fessler:18Asilomar}.\n\n\n\n\nThe remainder of this paper is organized as follows. \nSection~\\ref{sec:back} reviews how learned regularizers can help solve inverse problems.\nSection~\\ref{sec:CAOL} proposes the two CAOL models.\nSection~\\ref{sec:reBPG-M} introduces BPEG-M with several generalizations, analyzes its convergence, and applies a momentum coefficient formula and restarting technique from \\cite{Chun&Fessler:18TIP}.\nSection~\\ref{sec:CAOL+BPGM} applies the proposed BPEG-M methods to the CAOL models, designs two majorization matrices, and describes memory flexibility and applicability of parallel computing to BPEG-M-based CAOL.\nSection~\\ref{sec:CTrecon} introduces the CT MBIR model using a convolutional regularizer learned via CAOL \\cite{Chun&Fessler:18Asilomar}, along with its properties, i.e., its mathematical relation to a convolutional autoencoder, the importance of TF filters, and its algorithmic role in signal recovery.\nSection~\\ref{sec:result} reports numerical experiments that show \\textit{1)} the importance of sharp majorization in accelerating BPEG-M, and \\textit{2)} the benefits of BPEG-M-based CAOL -- acceleration, convergence, and memory flexibility.\nAdditionally, Section~\\ref{sec:result} reports sparse-view CT experiments that show \\textit{3)} the CT MBIR using learned convolutional regularizers significantly improves the reconstruction quality compared to that using a conventional edge-preserving (EP) regularizer, and \\textit{4)} more and wider filters in a learned regularizer better preserves edges in reconstructed images.\nFinally, Appendix~\\ref{sec:CNN} mathematically formulates unsupervised training of CNNs via CAOL, and shows that its updates attained via BPEG-M correspond to the three important CNN operators.\nAppendix~\\ref{sec:egs} introduces some potential applications of CAOL to image processing, imaging, and computer vision.\n\n\n\n\n\n\n\n\\section{Backgrounds: MBIR Using \\emph{Learned} Regularizers} \\label{sec:back}\n\nTo recover a signal $x \\in \\bbC^{N'}$ from a data vector $y \\in \\bbC^m$, \none often considers the following MBIR optimization problem\n(Appendix~\\ref{sec:notation} provides mathematical notations): \n$\n\\argmin_{x \\in \\cX} f(x; y) + \\gamma \\, g(x),\n$\nwhere $\\cX$ is a feasible set,\n$f(x; y)$ is data fidelity function that models imaging physics (or image formation) and noise statistics, \n$\\gamma \\!>\\! 0$ is a regularization parameter,\nand $g(x)$ is a regularizer, such as total variation \\cite[\\S2--3]{Arridge&etal:19AN}. \nHowever, when inverse problems are extremely ill-conditioned, \nthe MBIR approach using hand-crafted regularizers $g(x)$ has limitations in recovering signals.\nAlternatively, there has been a growing trend in learning sparsifying regularizers\n(e.g., convolutional regularizers \\cite{Chun&Fessler:18TIP, Chun&Fessler:17SAMPTA, Chun&etal:19SPL, Chun&Fessler:18Asilomar, Crockett&etal:19CAMSAP}) \nfrom training datasets and applying the learned regularizers to the following MBIR problem \\cite{Arridge&etal:19AN}:\n\\be{\n\\label{eq:mbir:learn}\n\\argmin_{x \\in \\cX} f(x; y) + \\gamma g(x; \\cO^\\star),\n\\tag{B1}\n}\nwhere a learned regularizer $g(x; \\cO^\\star)$ quantifies consistency between \nany candidate $x$ and training data that is encapsulated in some trained sparsifying operators $\\cO^\\star$.\nThe diagram in Fig.~\\ref{diag:abstract} shows the general process from\ntraining sparsifying operators to solving inverse problems via \\R{eq:mbir:learn}.\nSuch models \\R{eq:mbir:learn} arise in a wide range of applications. \nSee some examples in Appendix~\\ref{sec:egs}.\n\nThis paper describes multiple aspects of learning convolutional regularizers. \nThe next section first starts with proposing a new convolutional regularizer.\n\n\n\n\n\n\\section{CAOL: Models \\textit{Learning} Convolutional Regularizers} \\label{sec:CAOL}\n\nThe goal of CAOL is to find a set of filters that \\dquotes{best} sparsify a set of training images.\nCompared to hand-crafted regularizers,\nlearned convolutional regularizers \ncan better extract \\dquotes{true} features of estimated images and \nremove \\dquotes{noisy} features with thresholding operators.\nWe propose the following CAOL model:\n\\ea{\n\\label{sys:CAOL}\n\\argmin_{D = [d_1, \\ldots, d_K]} \\min_{\\{ z_{l,k} \\}} &~ F(D, \\{ z_{l,k} \\}) + \\beta g (D), \\tag{P0}\n\\\\\nF(D, \\{ z_{l,k} \\}) & := \\sum_{l=1}^L \\sum_{k=1}^K \\frac{1}{2} \\left\\| d_k \\circledast x_l - z_{l,k} \\right\\|_2^2 + \\alpha \\| z_{l,k} \\|_0, \\nn\n} \nwhere $\\circledast$ denotes a convolution operator (see details about boundary conditions in the supplementary material), $\\{ x_l \\in \\bbC^N : l =1,\\ldots,L \\}$ is a set of training images, $\\{ d_{k} \\in \\bbC^{R}: k = 1,\\ldots, K \\}$ is a set of convolutional kernels, $\\{ z_{l,k} \\in \\bbC^N : l = 1,\\ldots,L, k=1,\\ldots,K \\}$ is a set of sparse codes, and $g ( D )$ is a regularizer or constraint that encourages filter diversity or incoherence, $\\alpha \\!>\\! 0$ is a thresholding parameter controlling the sparsity of features $\\{ z_{l,k} \\}$, and $\\beta > 0$ is a regularization parameter for $g (D)$.\nWe group the $K$ filters into a matrix $D \\in \\bbC^{R \\times K}$:\n\\be{\n\\label{eq:D}\nD := \\left[ \\arraycolsep=3pt \\begin{array}{ccc} d_1 & \\ldots & d_K \\end{array} \\right].\n}\nFor simplicity, we fix the dimension for training signals, i.e., $\\{ x_l, z_{l,k} \\in \\bbC^N \\}$, but the proposed model \\R{sys:CAOL} can use training signals of different dimension, i.e., $\\{ x_l, z_{l,k} \\in \\bbC^{N_l} \\}$.\nFor sparse-view CT in particular, \nthe diagram in Fig.~\\ref{diag:caol-ctmbir} shows the process from\nCAOL \\R{sys:CAOL} to solving its inverse problem via\nMBIR using learned convolutional regularizers.\n\nThe following two subsections design the constraint or regularizer $g(D)$ \nto avoid redundant filters (without it, all filters could be identical).\n\n\n\n\n\n\n\\begin{figure*}[!pt]\n\\centering\n\n\\begin{tabular}{c}\n\\includegraphics[trim={2.2cm 9.5cm 2.2cm 9.9cm},clip,scale=0.46]{.\/Fig_v1\/diagram_caol-ctmbir.pdf}\n\\end{tabular}\n\n\\vspace{-0.25em}\n\\caption{A flowchart from CAOL \\R{sys:CAOL} to \nMBIR using a convolutional sparsifying regularizer learned via CAOL \\R{sys:CT&CAOL}\nin sparse-view CT.\nSee details of the CAOL process \\R{sys:CAOL} and its variants \\R{sys:CAOL:orth}--\\R{sys:CAOL:div}, \nand the CT MBIR process \\R{sys:CT&CAOL} in \nSection~\\ref{sec:CAOL} and Section~\\ref{sec:CTrecon}, respectively.\n}\n\\label{diag:caol-ctmbir}\n\\end{figure*}\n\n\n\n\n\n\n\\subsection{CAOL with Orthogonality Constraint} \\label{sec:CAOL:orth}\n\nWe first propose a CAOL model with a nonconvex orthogonality constraint on the filter matrix $D$ in \\R{eq:D}:\n\\be{\n\\label{sys:CAOL:orth}\n\\argmin_{D} \\min_{\\{ z_{l,k} \\}} ~ F(D, \\{ z_{l,k} \\}) \\quad \\mathrm{subj.~to} ~ D D^H = \\frac{1}{R} \\cdot I. \\tag{P1}\n} \nThe orthogonality condition $D D^H = \\frac{1}{R} I$ in \\R{sys:CAOL:orth} enforces a TF condition on the filters $\\{ d_k \\}$ in CAOL \\R{sys:CAOL}.\nProposition~\\ref{p:TFconst} below formally states this relation.\n\n\\prop{[Tight-frame filters]\\label{p:TFconst}\nFilters satisfying the orthogonality constraint \n$D D^H = \\frac{1}{R} I$ in \\R{sys:CAOL:orth} satisfy the following TF condition in a convolution perspective:\n\\be{\n\\label{eq:CAOL:TFcond}\n\\sum_{k=1}^K \\nm{ d_k \\circledast x }_2^2 = \\nm{ x }_2^2, \\quad \\forall x \\in \\bbC^N,\n}\nfor both circular and symmetric boundary conditions. \n}\n\n\\prf{\\renewcommand{\\qedsymbol}{}\nSee Section~\\ref{sec:prf:p:TF} of the supplementary material.\n}\n\n\nProposition~\\ref{p:TFconst} corresponds to a TF result from patch-domain approaches; see~Section~\\ref{sec:prf:p:TF}. (Note that the patch-domain approach in \\cite[Prop.~3]{Cai&etal:14ACHA} requires $R = K$.)\nHowever, we constrain the filter dimension to be $R \\leq K$ to have an efficient solution for CAOL model \\R{sys:CAOL:orth}; see Proposition~\\ref{p:orth} later.\nThe following section proposes a more flexible CAOL model in terms of the filter dimensions $R$ and $K$.\n\n\n\n\n\\subsection{CAOL with Diversity Promoting Regularizer}\n\nAs an alternative to the CAOL model \\R{sys:CAOL:orth}, we propose a CAOL model with a diversity promoting regularizer and a nonconvex norm constraint on the filters $\\{ d_k \\}$:\n\\ea{\n\\label{sys:CAOL:div}\n\\argmin_{D} \\min_{\\{ z_{l,k} \\}} &~ F(D, \\{ z_{l,k} \\}) + \\frac{\\beta}{2} \\overbrace{\\left\\| D^H D - \\frac{1}{R} \\cdot I \\right\\|_{\\mathrm{F}}^2}^{\\mathrm{\\hbox{$=: g_{\\text{div}} (D)$}}}, \\nn\n\\\\\n\\mathrm{subject~to} ~&~ \\| d_{k} \\|_2^2 = \\frac{1}{R}, \\quad k = 1,\\ldots, K. \\tag{P2}\n} \nIn the CAOL model \\R{sys:CAOL:div}, we consider the following:\n\\bulls{\n\\item The constraint in \\R{sys:CAOL:div} forces the learned filters $\\{ d_k \\}$ to have uniform energy. In addition, it avoids the \\dquotes{scale ambiguity} problem \\cite{Rem&Karin:10TIT}. \n\n\\item The regularizer in \\R{sys:CAOL:div}, $g_{\\text{div}} (D)$, promotes filter diversity, i.e., incoherence between $d_k$ and $\\{ d_{k'} : k' \\neq k \\}$, measured by $| \\ip{d_k}{d_{k'}} |^2$ for $k \\neq k'$.\n}\n\nWhen $R = K$ and $\\beta \\rightarrow \\infty$, the model \\R{sys:CAOL:div} becomes \\R{sys:CAOL:orth} since $D^H D = \\frac{1}{R} I$ implies $D D^H = \\frac{1}{R} I$ (for square matrices $A$ and $B$, if $AB=I$ then $BA = I$).\nThus \\R{sys:CAOL:div} generalizes \\R{sys:CAOL:orth} by relaxing the off-diagonal elements of the equality constraint in \\R{sys:CAOL:orth}. (In other words, when $R=K$, the orthogonality constraint in \\R{sys:CAOL:orth} enforces the TF condition and promotes the filter diversity.)\nOne price of this generalization is the extra tuning parameter $\\beta$.\n\n\n\n\\R{sys:CAOL:orth}--\\R{sys:CAOL:div} are challenging nonconvex optimization problems and block optimization approaches seem suitable. \nThe following section proposes a new block optimization method with momentum and majorizers, \nto rapidly solve the multiple block multi-nonconvex problems proposed in this paper, while guaranteeing convergence to critical points.\n\n\n\n\n\n\n\n\n\n\\section{BPEG-M: Solving Block Multi-Nonconvex Problems with Convergence Guarantees} \\label{sec:reBPG-M}\n\nThis section describes a new optimization approach, BPEG-M,\nfor solving block multi-nonconvex problems like \n\\textit{a)} CAOL \\R{sys:CAOL:orth}--\\R{sys:CAOL:div},\\footnote{\nA block coordinate descent algorithm can be applied to CAOL \\R{sys:CAOL:orth};\nhowever, its convergence guarantee in solving CAOL \\R{sys:CAOL:orth} is not yet known and might require stronger sufficient conditions than BPEG-M \\cite{Tseng:2001JOTA}.\n} \n\\textit{b)} CT MBIR \\R{sys:CT&CAOL} using learned convolutional regularizer via \\R{sys:CAOL:orth} (see Section~\\ref{sec:CTrecon}),\nand \\textit{c)} \\dquotes{hierarchical} CAOL \\R{sys:CNN:orth} (see Appendix~\\ref{sec:CNN}).\n\n\n\\subsection{BPEG-M -- Setup} \\label{sec:BPGM:setup}\n\nWe treat the variables of the underlying optimization problem either as a single block or multiple disjoint blocks. \nSpecifically, consider the following \\textit{block multi-nonconvex} optimization problem:\n\\ea{\n\\label{sys:multiConvx}\n\\min &~ F(x_1,\\ldots, x_B) := f(x_1,\\ldots,x_B) + \\sum_{b=1}^B g_b (x_b),\n}\nwhere variable $x$ is decomposed into $B$ blocks $x_1 ,\\ldots,x_B$ ($\\{ x_b \\in \\bbR^{n_b} : b=1,\\ldots,B \\}$), $f$ is assumed to be continuously differentiable, but functions $\\{ g_b : b = 1,\\ldots,B \\}$ are not necessarily differentiable. \nThe function $g_b$ can incorporate the constraint $x_b \\in \\cX_b$, by allowing any $g_b$ to be extended-valued, e.g., $g_b (x_b) = \\infty$ if $x_b \\notin \\cX_b$, for $b=1,\\ldots,B$.\nIt is standard to assume that both $f$ and $\\{ g_b \\}$ are closed and proper and the sets $\\{ \\cX_b \\}$ are closed and nonempty.\nWe do \\textit{not} assume that $f$, $\\{ g_b \\}$, or $\\{ \\cX_b \\}$ are convex. Importantly, $g_b$ can be a nonconvex $\\ell^p$ quasi-norm, $p \\in [0, 1)$.\nThe general block multi-convex problem in \\cite{Chun&Fessler:18TIP, Xu&Yin:13SIAM} is a special case of \\R{sys:multiConvx}.\n\n\nThe BPEG-M framework considers a more general concept than Lipschitz continuity of the gradient as follows:\n\n\\defn{[$M$-Lipschitz continuity] \\label{d:QM}\nA function $g: \\bbR^n \\rightarrow \\bbR^{n}$ is \\emph{$M$-Lipschitz continuous} on $\\bbR^n$ if there exist a (symmetric) positive definite matrix $M$ such that\n\\bes{\n\\nm{g(x) - g(y)}_{M^{-1}} \\leq \\nm{x - y}_{M}, \\quad \\forall x,y,\n}\nwhere $\\nm{x}_{M}^2 := x^T M x$.\n}\n\nLipschitz continuity is a special case of $M$-Lipschitz continuity with $M$ equal to a scaled identity matrix with \na Lipschitz constant of the gradient $\\nabla f$ (e.g., for $f(x) = \\frac{1}{2} \\| Ax - b\\|_2^2$, the (smallest) Lipschitz constant of $\\nabla f$ is the maximum eigenvalue of $A^T A$).\nIf the gradient of a function is $M$-Lipschitz continuous, then we obtain the following quadratic majorizer (i.e., surrogate function \\cite{Lange&Hunter&Yang:00JCGS, Jacobson&Fessler:07TIP}) at a given point $y$ without assuming convexity:\n\n\\lem{[Quadratic majorization (QM) via $M$-Lipschitz continuous gradients]\n\\label{l:QM}\nLet $f : \\bbR^n \\rightarrow \\bbR$. If $\\nabla f$ is $M$-Lipschitz continuous, then\n\\bes{\nf(x) \\leq f(y) + \\ip{\\nabla f(y)}{x-y} + \\frac{1}{2} \\nm{x - y}_M^2, \\quad \\forall x,y \\in \\bbR^n.\n}\n}\n\n\\prf{\\renewcommand{\\qedsymbol}{}\nSee Section~\\ref{sec:prf:l:QM} of the supplementary material.\n}\n\nExploiting Definition~\\ref{d:QM} and Lemma~\\ref{l:QM}, the proposed method, BPEG-M, is given as follows.\nTo solve \\R{sys:multiConvx}, we minimize a majorizer of $F$ cyclically over each block $x_1,\\ldots,x_B$, while fixing the remaining blocks at their previously updated variables. Let $x_b^{(i+1)}$ be the value of $x_b$ after its $i\\rth$ update, and define\n\\begingroup\n\\setlength{\\thinmuskip}{1.5mu}\n\\setlength{\\medmuskip}{2mu plus 1mu minus 2mu}\n\\setlength{\\thickmuskip}{2.5mu plus 2.5mu}\n\\bes{\nf_b^{(i+1)}(x_b) := f \\Big( x_1^{(i+1)}, \\ldots, x_{b-1}^{(i+1)}, x_b, x_{b+1}^{(i)}, \\ldots, x_{B}^{(i)} \\Big), \\quad \\forall b,i.\n}\n\\endgroup\nAt the $b\\rth$ block of the $i\\rth$ iteration, we apply Lemma~\\ref{l:QM} to functional $f_b^{(i+1)}(x_b)$ with a $M_b^{(i+1)}$-Lipschitz continuous gradient, and minimize the majorized function.\\footnote{The quadratically majorized function allows a unique minimizer if $g_b^{(i+1)} (x_b)$ is convex and $\\cX_b^{(i+1)}$ is a convex set (note that $M_b^{(i+1)} \\!\\succ\\! 0$).} Specifically, BPEG-M uses the updates\n\\begingroup\n\\setlength{\\thinmuskip}{1.5mu}\n\\setlength{\\medmuskip}{2mu plus 1mu minus 2mu}\n\\setlength{\\thickmuskip}{2.5mu plus 2.5mu}\n\\fontsize{9.5pt}{11.4pt}\\selectfont\n\\allowdisplaybreaks\n\\ea{\nx_b^{(i+1)} \n& = \\argmin_{ x_b } \\hspace{0.1em} \\ip{ \\nabla_{x_b} f_b^{(i+1)} (\\acute{x}_b^{(i+1)}) }{ x_b - \\acute{x}_b^{(i+1)} } \\nn\n\\\\\n& \\hspace{4.5em} + \\frac{1}{2} \\nm{ x_b - \\acute{x}_b^{(i+1)} }_{\\widetilde{M}_b^{(i+1)}}^2 + g_b (x_b) \\label{update:x} \n\\\\\n&= \\argmin_{ x_b } \\hspace{0.1em} \\frac{1}{2} \\bigg\\| x_b - \\bigg( \\acute{x}_b^{(i+1)} - \\Big( \\widetilde{M}_b^{(i+1)} \\Big)^{\\!\\!\\!-1} \\nn\n\\\\\n& \\hspace{8.0em} \\cdot \\nabla_{x_b} f_b^{(i+1)} (\\acute{x}_b^{(i+1)}) \\bigg) \\bigg\\|_{\\widetilde{M}_b^{(i+1)}}^2 \\!\\! + g_b (x_b) \\nn\n\\\\\n&= \\mathrm{Prox}_{g_b}^{\\widetilde{M}_b^{(i+1)}} \\!\\! \\bigg( \\! \\underbrace{ \\acute{x}_b^{(i+1)} - \\Big( \\! \\widetilde{M}_b^{(i+1)} \\! \\Big)^{\\!\\!\\!-1} \\! \\nabla_{x_b} f_b^{(i+1)} (\\acute{x}_b^{(i+1)}) }_{\\text{\\emph{extrapolated gradient step using a majorizer} of $f_b^{(i+1)}$}} \\! \\bigg), \\nn\n}\n\\endgroup\nwhere\n\\be{\n\\label{update:xacute}\n\\acute{x}_b^{(i+1)} = x_{b}^{(i)} + E_b^{(i+1)} \\left( x_b^{(i)} - x_b^{(i-1)} \\right),\n}\nthe proximal operator is defined by\n\\bes{\n\\mathrm{Prox}_g^M (y) := \\argmin_{x} \\, \\frac{1}{2} \\nm{x - y}_M^2 + g(x),\n}\n$\\nabla f_b^{(i+1)} (\\acute{x}_b^{(i+1)})$ is the block-partial gradient of $f$ at $\\acute{x}_b^{(i+1)}$, an \\textit{upper-bounded majorization matrix} is updated by\n\\be{\n\\label{update:Mtilde}\n\\widetilde{M}_b^{(i+1)} = \\lambda_b \\cdot M_b^{(i+1)} \\succ 0, \\qquad \\lambda_b > 1,\n}\nand $M_b^{(i+1)} \\!\\in\\! \\bbR^{n_b \\times n_b}$ is a symmetric positive definite \\textit{majorization matrix} of $\\nabla f_b^{(i+1)}$.\nIn \\R{update:xacute}, the $\\bbR^{n_b \\times n_b}$ matrix $E_b^{(i+1)} \\succeq 0$ is an \\textit{extrapolation matrix} that accelerates convergence in solving block multi-convex problems \\cite{Chun&Fessler:18TIP}.\nWe design it in the following form:\n\\be{\n\\label{update:Eb}\nE_b^{(i+1)} = e_b^{(i)} \\cdot \\frac{\\delta (\\lambda_b - 1)}{2 (\\lambda_b + 1)} \\cdot \\left( M_b^{(i+1)} \\right)^{-1\/2} \\left( M_b^{(i)} \\right)^{1\/2}, \n}\nfor some $\\{ 0 \\leq e_b^{(i)} \\leq 1 : \\forall b, i \\}$ and $\\delta < 1$, to satisfy condition \\R{cond:Wb} below.\nIn general, choosing $\\lambda_b$ values in \\R{update:Mtilde}--\\R{update:Eb} to accelerate convergence is application-specific.\nAlgorithm~\\ref{alg:BPGM} summarizes these updates.\n\n\n\n\n\n\n\\begin{algorithm}[t!]\n\\caption{BPEG-M}\n\\label{alg:BPGM}\n\n\\begin{algorithmic}\n\\REQUIRE $\\{ x_b^{(0)} = x_b^{(-1)} : \\forall b \\}$, $\\{ E_b^{(i)} \\in [0, 1], \\forall b,i \\}$, $i=0$\n\n\\WHILE{a stopping criterion is not satisfied}\n\n\\FOR{$b = 1,\\ldots,B$}\n\n\\begingroup\n\\setlength{\\thinmuskip}{1.5mu}\n\\setlength{\\medmuskip}{2mu plus 1mu minus 2mu}\n\\setlength{\\thickmuskip}{2.5mu plus 2.5mu}\n\\fontsize{9.5pt}{11.4pt}\\selectfont\n\\STATE Calculate $M_b^{(i+1)}$, $\\displaystyle \\widetilde{M}_b^{(i+1)}$ by \\R{update:Mtilde}, and $E_b^{(i+1)}$ by \\R{update:Eb}\n\\STATE $\\displaystyle \\acute{x}_b^{(i+1)} = \\, x_{b}^{(i)} + E_b^{(i+1)} \\! \\left( x_b^{(i)} - x_b^{(i-1)} \\right)$\n\\STATE $\\displaystyle x_b^{(i+1)} = \\, \\ldots$ \\\\ $\\displaystyle \\mathrm{Prox}_{g_b}^{\\widetilde{M}_b^{(i+1)}} \\!\\! \\bigg( \\!\\! \\acute{x}_b^{(i+1)} \\!-\\! \\Big( \\! \\widetilde{M}_b^{(i+1)} \\! \\Big)^{\\!\\!\\!\\!-1} \\!\\! \\nabla f_b^{(i+1)} (\\acute{x}_b^{(i+1)}) \\!\\! \\bigg)$\n\\endgroup\n\n\\ENDFOR\n\n\\STATE $i = i+1$\n\n\\ENDWHILE\n\n\\end{algorithmic}\n\\end{algorithm}\n\n\n\n\n\n\nThe majorization matrices $M_b^{(i)}$ and $\\widetilde{M}_b^{(i+1)}$ in \\R{update:Mtilde} influence the convergence rate of BPEG-M.\nA tighter majorization matrix (i.e., a matrix giving tighter bounds in the sense of Lemma~\\ref{l:QM}) provided faster convergence rate \\cite[Lem.~1]{Fessler&etal:93TNS}, \\cite[Fig.~2--3]{Chun&Fessler:18TIP}.\nAn interesting observation in Algorithm~\\ref{alg:BPGM} is that there exists a tradeoff between majorization sharpness via \\R{update:Mtilde} and extrapolation effect via \\R{update:xacute} and \\R{update:Eb}. For example, increasing $\\lambda_b$ (e.g., $\\lambda_b = 2$) allows more extrapolation but results in looser majorization; setting $\\lambda_b \\rightarrow 1$ results in sharper majorization but provides less extrapolation.\n\n\n\\rem{\\label{r:BPGM}\nThe proposed BPEG-M framework -- with key updates \\R{update:x}--\\R{update:xacute} -- generalizes the BPG method \\cite{Xu&Yin:17JSC}, and has several benefits over BPG \\cite{Xu&Yin:17JSC} and BPEG-M introduced earlier in \\cite{Chun&Fessler:18TIP}:\n\\begin{itemize}\n\\item The BPG setup in \\cite{Xu&Yin:17JSC} is a particular case of BPEG-M using a scaled identity majorization matrix $M_b$ with a Lipschitz constant of $\\nabla f_b^{(i+1)} (\\acute{x}_b^{(i+1)})$. The BPEG-M framework can significantly accelerate convergence by allowing sharp majorization; see \\cite[Fig.~2--3]{Chun&Fessler:18TIP} and Fig.~\\ref{fig:Comp:diffBPGM}. \nThis generalization was first introduced for block multi-convex problems in \\cite{Chun&Fessler:18TIP}, \nbut the proposed BPEG-M in this paper addresses the more general problem, block multi-(non)convex optimization.\n\n\\item BPEG-M is useful for controlling the tradeoff between majorization sharpness and extrapolation effect in different blocks, by allowing each block to use different $\\lambda_b$ values. If tight majorization matrices can be designed for a certain block $b$, then it could be reasonable to maintain the majorization sharpness by setting $\\lambda_b$ very close to 1. \nWhen setting $\\lambda_b = 1 + \\epsilon$ (e.g., $\\epsilon$ is a machine epsilon) and using $E_b^{(i+1)} = 0$ (no extrapolation), solutions of the original and its upper-bounded problem become (almost) identical. In such cases, it is unnecessary to solve the upper bounded problem \\R{update:x}, and the proposed BPEG-M framework allows using the solution of $f_b^{(i+1)}(x_b)$ without QM; see Section~\\ref{sec:CAOL:spCd}. This generalization was not considered in \\cite{Xu&Yin:17JSC}.\n\n\\item The condition for designing the extrapolation matrix \\R{update:Eb}, i.e., \\R{cond:Wb} in Assumption 3, is more general than that in \\cite[(9)]{Chun&Fessler:18TIP} (e.g., \\R{eg:Wb:eigen}).\nSpecifically, the matrices $E_b^{(i+1)}$ and $M_b^{(i+1)}$ in \\R{update:Eb} need not be diagonalized by the same basis.\n\\end{itemize}\n}\n\n\n\n\n\n\n\nThe first two generalizations lead to the question, \\dquotes{Under the sharp QM regime (i.e., having tight bounds in Lemma~\\ref{l:QM}), what is the best way in controlling $\\{ \\lambda_b \\}$ in \\R{update:Mtilde}--\\R{update:Eb} in Algorithm~\\ref{alg:BPGM}?}\nOur experiments show that, if sufficiently sharp majorizers are obtained for partial or all blocks, then giving more weight to sharp majorization provides faster convergence compared to emphasizing extrapolation; for example, $\\lambda_{b} = 1+\\epsilon$ gives faster convergence than $\\lambda_b = 2$.\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{BPEG-M -- Convergence Analysis} \\label{sec:convg-analysis}\n\nThis section analyzes the convergence of Algorithm~\\ref{alg:BPGM} under the following assumptions.\n\n\\begin{itemize}\n\\item[] {\\em Assumption 1)} $F$ is proper and lower bounded in $\\dom(F)$, $f$ is continuously differentiable, $g_b$ is proper lower semicontinuous, $\\forall b$.\\footnote{\n$F : \\bbR^n \\rightarrow (- \\infty, + \\infty]$ is proper if $\\dom F \\neq \\emptyset$. \n$F$ is lower bounded in $\\dom(F) := \\{ x : F(x) < \\infty \\}$ if $\\inf_{x \\in \\dom(F)} F(x) > -\\infty$.\n$F$ is lower semicontinuous at point $x_0$ if $\\liminf_{x \\rightarrow x_0} F(x) \\geq F(x_0)$.\n}\n\\R{sys:multiConvx} has a critical point $\\bar{x}$, i.e., $0 \\in \\partial F(\\bar{x})$,\nwhere $\\partial F(x)$ denotes the limiting subdifferential of $F$ at $x$ (see \\cite[\\S 1.9]{Kruger:03JMS}, \\cite[\\S 8]{Rockafellar&Wets:book}).\n\\item[] {\\em Assumption 2)} The block-partial gradients of $f$, $\\nabla f_b^{(i+1)}$, are $M_b^{(i+1)}$-Lipschitz continuous, i.e.,\n\\ea{\n\\label{eq:QMbound}\n&~ \\nm{ \\nabla_{x_b} f_b^{(i+1)} (u) - \\nabla_{x_b} f_b^{(i+1)} (v) }_{\\big( \\! M_b^{(i+1)} \\! \\big)^{\\!-1}} \n\\nn\n\\\\\n& \\leq \\nm{u - v}_{M_b^{(i+1)}},\n}\nfor $u,v \\in \\bbR^{n_b}$, and (unscaled) majorization matrices satisfy $m_{b} I_{n_b} \\preceq M_b^{(i+1)}$ with $0 < m_{b} < \\infty$, $\\forall b, i$.\n\\item[] {\\em Assumption 3)} The extrapolation matrices $E_b^{(i+1)} \\succeq 0$ satisfy\n\\be{\n\\label{cond:Wb}\n\\Big( E_b^{(i+1)} \\Big)^T M_b^{(i+1)} E_b^{(i+1)} \\preceq \\frac{\\delta^2 (\\lambda_b - 1)^2}{4 (\\lambda_b + 1)^2} \\cdot M_b^{(i)},\n}\nfor any $\\delta < 1$, $\\forall b,i$.\n\\end{itemize}\n\n\nCondition \\R{cond:Wb} in Assumption 3 generalizes that in \\cite[Assumption~3]{Chun&Fessler:18TIP}.\nIf eigenspaces of $E_b^{(i+1)}$ and $M_b^{(i+1)}$ coincide (e.g., diagonal and circulant matrices), $\\forall i$ \\cite[Assumption~3]{Chun&Fessler:18TIP}, \\R{cond:Wb} becomes\n\\begingroup\n\\setlength{\\thinmuskip}{1.5mu}\n\\setlength{\\medmuskip}{2mu plus 1mu minus 2mu}\n\\setlength{\\thickmuskip}{2.5mu plus 2.5mu}\n\\be{\n\\label{eg:Wb:eigen}\nE_b^{(i+1)} \\preceq \\frac{\\delta (\\lambda_b - 1)}{2 (\\lambda_b + 1)} \\cdot \\left( M_b^{(i)} \\right)^{1\/2} \\left( M_b^{(i+1)} \\right)^{-1\/2},\n} \n\\endgroup\nas similarly given in \\cite[(9)]{Chun&Fessler:18TIP}. \nThis generalization allows one to consider arbitrary structures of $M_b^{(i)}$ across iterations. \n\n\n\n\\lem{[Sequence bounds]\n\\label{l:seqBound}\nLet $\\{ \\widetilde{M}_b : b = 1,\\ldots,B \\}$ and $\\{ E_b : b = 1,\\ldots,B \\}$ be as in \\R{update:Mtilde}--\\R{update:Eb}, respectively. \nThe cost function decrease for the $i\\rth$ update satisfies:\n\\ea{\n\\label{eq:l:seqBound}\nF_b (x_b^{(i)}) - F_b (x_b^{(i+1)}) \n& \\geq \\frac{\\lambda_b - 1}{4} \\nm{ x_b^{(i)} - x_b^{(i+1)} }_{M_b^{(i+1)}}^2 \\nn\n\\\\\n& \\hspace{1.2em} - \\frac{ ( \\lambda_b - 1 ) \\delta^2}{4} \\nm{ x_b^{(i-1)} - x_b^{(i)} }_{ M_b^{(i)}}^2\n}\n}\n\\prf{\\renewcommand{\\qedsymbol}{}\nSee Section~\\ref{sec:prf:l:seqBound} of the supplementary material.\n}\n\n \nLemma~\\ref{l:seqBound} generalizes \\cite[Lem.~1]{Xu&Yin:17JSC} using $\\{ \\lambda_b = 2 \\}$. \nTaking the majorization matrices in \\R{eq:l:seqBound} to be scaled identities with Lipschitz constants, i.e., $M_b^{(i+1)} \\!=\\! L_b^{(i+1)} \\cdot I$ and $M_b^{(i)} \\!=\\! L_b^{(i)} \\cdot I$, where $L_b^{(i+1)}$ and $L_b^{(i)}$ are Lipschitz constants, the bound \\R{eq:l:seqBound} becomes equivalent to that in \\cite[(13)]{Xu&Yin:17JSC}.\nNote that BPEG-M for block multi-convex problems in \\cite{Chun&Fessler:18TIP} can be viewed within BPEG-M in Algorithm~\\ref{alg:BPGM}, by similar reasons in \\cite[Rem.~2]{Xu&Yin:17JSC} -- bound \\R{eq:l:seqBound} holds for the block multi-convex problems by taking $E_b^{(i+1)}$ in \\R{eg:Wb:eigen} as $E_b^{(i+1)} \\preceq \\delta \\cdot ( M_b^{(i)} )^{1\/2} ( M_b^{(i+1)} )^{-1\/2}$ in \\cite[Prop.~3.2]{Chun&Fessler:18TIP}. \n\n\n\n\n\n\n\\prop{[Square summability]\n\\label{p:sqSum}\nLet $\\{ x^{(i+1)} : i \\geq 0 \\}$ be generated by Algorithm \\ref{alg:BPGM}.\nWe have\n\\be{\n\\label{eq:p:sqSum}\n\\sum_{i=0}^{\\infty} \\nm{ x^{(i)} - x^{(i+1)} }_2^2 < \\infty.\n}\n}\n\n\\prf{\\renewcommand{\\qedsymbol}{}\nSee Section~\\ref{sec:prf:p:sqSum} of the supplementary material.\n}\n\n\n\n\n\nProposition~\\ref{p:sqSum} implies that\n\\be{\n\\label{eq:convg_to0}\n\\nm{ x^{(i)} - x^{(i+1)} }_2^2 \\rightarrow 0,\n}\nand \\R{eq:convg_to0} is used to prove the following theorem:\n\n\n\n\n\n\\thm{[A limit point is a critical point]\n\\label{t:subseqConv}\nUnder Assumptions 1--3, let $\\{ x^{(i+1)} : i \\geq 0 \\}$ be generated by Algorithm \\ref{alg:BPGM}.\nThen any limit point $\\bar{x}$ of $\\{ x^{(i+1)} : i \\geq 0 \\}$ is a critical point of \\R{sys:multiConvx}. \nIf the subsequence $\\{ x^{(i_j+1)} \\}$ converges to $\\bar{x}$, then\n\\bes{\n\\lim_{j \\rightarrow \\infty} F(x^{(i_j + 1 )}) = F(\\bar{x}).\n} \n}\n\n\\prf{\\renewcommand{\\qedsymbol}{}\nSee Section~\\ref{sec:prf:t:subseqConv} of the supplementary material.\n}\n\nFinite limit points exist if the generated sequence $\\{ x^{(i+1)} : i \\geq 0 \\}$ is bounded; see, for example, \\cite[Lem.~3.2--3.3]{Bao&Ji&Shen:14ACHA}. For some applications, the boundedness of $\\{ x^{(i+1)} : i \\geq 0 \\}$ can be satisfied by choosing appropriate regularization parameters, e.g., \\cite{Chun&Fessler:18TIP}.\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{Restarting BPEG-M} \\label{sec:reBPGM}\n\nBPG-type methods \\cite{Chun&Fessler:18TIP, Xu&Yin:17JSC, Xu&Yin:13SIAM} can be further accelerated by applying \\textit{1)} a momentum coefficient formula similar to those used in fast proximal gradient (FPG) methods \\cite{Beck&Teboulle:09SIAM, Nesterov:07CORE, Tseng:08techRep}, and\/or \\textit{2)} an adaptive momentum restarting scheme \\cite{ODonoghue&Candes:15FCM, Giselsson&Boyd:14CDC}; see \\cite{Chun&Fessler:18TIP}. \nThis section applies these two techniques to further accelerate BPEG-M in Algorithm~\\ref{alg:BPGM}.\n\nFirst, we apply the following increasing momentum-coefficient formula to \\R{update:Eb} \\cite{Beck&Teboulle:09SIAM}: \n\\be{\n\\label{eq:mom_coeff}\ne_b^{(i+1)} = \\frac{\\theta^{(i)} - 1}{\\theta^{(i+1)}}, \\quad \\theta^{(i+1)} = \\frac{1 + \\sqrt{1 + 4 (\\theta^{(i)})^2}}{2}.\n}\nThis choice guarantees fast convergence of FPG method \\cite{Beck&Teboulle:09SIAM}.\nSecond, we apply a momentum restarting scheme \\cite{ODonoghue&Candes:15FCM, Giselsson&Boyd:14CDC}, when the following \\textit{gradient-mapping} criterion is met \\cite{Chun&Fessler:18TIP}:\n\\be{\n\\label{eq:restart:grad}\n\\cos \\! \\left( \\Theta \\! \\left( M_b^{(i+1)} \\!\\! \\left( \\acute{x}_b^{(i+1)} - x_b^{(i+1)} \\right), x_b^{(i+1)} - x_b^{(i)} \\right) \\right) > \\omega,\n}\nwhere the angle between two nonzero real vectors $\\vartheta$ and $\\vartheta'$ is $\\Theta (\\vartheta, \\vartheta') := \\ip{\\vartheta}{\\vartheta'} \/ ( \\nm{\\vartheta}_2 \\nm{\\vartheta'}_2 )$ and $\\omega \\in [-1, 0]$.\nThis scheme restarts the algorithm whenever the momentum, i.e., $x_b^{(i+1)} - x_b^{(i)}$, is likely to lead the algorithm in an unhelpful direction, as measured by the gradient mapping at the $x_b^{(i+1)}$-update.\nWe refer to BPEG-M combined with the methods \\R{eq:mom_coeff}--\\R{eq:restart:grad} as restarting BPEG-M (reBPEG-M). Section~\\ref{sec:reBPGM-supp} in the supplementary material summarizes the updates of reBPEG-M.\n\nTo solve the block multi-nonconvex problems proposed in this paper (e.g., \\R{sys:CAOL:orth}--\\R{sys:CT&CAOL}), we apply reBPEG-M (a variant of Algorithm~\\ref{alg:BPGM}; see Algorithm~\\ref{alg:reBPGM}), promoting fast convergence to a critical point. \n\n\n\n\n\n\n\n\n\n\n\\section{Fast and Convergent CAOL via BPEG-M} \\label{sec:CAOL+BPGM}\n\nThis section applies the general BPEG-M approach to CAOL.\nThe CAOL models \\R{sys:CAOL:orth} and \\R{sys:CAOL:div} satisfy the assumptions of BPEG-M; see Assumption 1--3 in Section~\\ref{sec:convg-analysis}.\nCAOL models \\R{sys:CAOL:orth} and \\R{sys:CAOL:div} readily satisfy Assumption~1 of BPEG-M.\nTo show the continuously differentiability of $f$ and the lower boundedness of $F$, consider that \\textit{1)} $\\sum_{l} \\sum_{k} \\frac{1}{2} \\left\\| d_k \\circledast x_l - z_{l,k} \\right\\|_2^2$ in \\R{sys:CAOL} is continuously differentiable with respect to $D$ and $\\{ z_{l,k} \\}$;\n\\textit{2)} the sequences $\\{ D^{(i+1)} \\}$ are bounded, because they are in the compact set $\\cD_{\\text{\\R{sys:CAOL:orth}}} = \\{ D: D D^H = \\frac{1}{R} I \\}$ and $\\cD_{\\text{\\R{sys:CAOL:div}}} = \\{ d_k : \\| d_k \\|_2^2 = \\frac{1}{R}, \\forall k \\}$ in \\R{sys:CAOL:orth} and \\R{sys:CAOL:div}, respectively; \nand \\textit{3)} the positive thresholding parameter $\\alpha$ ensures that the sequence $\\{ z_{l,k}^{(i+1)} \\}$ is bounded (otherwise the cost would diverge).\nIn addition, for both \\R{sys:CAOL:orth} and \\R{sys:CAOL:div}, the lower semicontinuity of regularizer $g_b$ holds, $\\forall b$. \nFor $D$-optimization, the indicator function of the sets $\\cD_{\\text{\\R{sys:CAOL:orth}}}$ and $\\cD_{\\text{\\R{sys:CAOL:div}}}$ is lower semicontinuous, because the sets are compact.\nFor $\\{z_{l,k}\\}$-optimization, the $\\ell^0$-quasi-norm is a lower semicontinuous function.\nAssumptions~2 and 3 are satisfied with the majorization matrix designs in this section -- see Sections~\\ref{sec:CAOL:filt}--\\ref{sec:CAOL:spCd} later -- and the extrapolation matrix design in \\R{update:Eb}, respectively.\n\nSince CAOL models \\R{sys:CAOL:orth} and \\R{sys:CAOL:div} satisfy the BPEG-M conditions, we solve \\R{sys:CAOL:orth} and \\R{sys:CAOL:div} by the reBPEG-M method with a two-block scheme, i.e., we alternatively update all filters $D$ and all sparse codes $\\{z_{l,k} : l = 1,\\ldots,L, k=1,\\ldots,K\\}$. \nSections~\\ref{sec:CAOL:filt} and \\ref{sec:CAOL:spCd} describe details of $D$-block and $\\{ z_{l,k} \\}$-block optimization within the BPEG-M framework, respectively.\nThe BPEG-M-based CAOL algorithm is particularly useful for learning convolutional regularizers from large datasets because of its memory flexibility and parallel computing applicability, as described in Section~\\ref{sec:CAOL:memory} and Sections~\\ref{sec:CAOL:filt}--\\ref{sec:CAOL:spCd}, respectively.\n\n\n\n\n\n\n\n\n\\subsection{Filter Update: $D$-Block Optimization} \\label{sec:CAOL:filt}\n\nWe first investigate the structure of the system matrix in the filter update for \\R{sys:CAOL}.\nThis is useful for \\textit{1)} accelerating majorization matrix computation in filter updates (e.g., Lemmas~\\ref{l:MDdiag}--\\ref{l:MDscaleI}) and \\textit{2)} applying $R \\!\\times\\! N$-sized adjoint operators (e.g., $\\Psi_l^H$ in \\R{eq:Psi_l} below) to an $N$-sized vector without needing the Fourier approach \\cite[Sec.~\\Romnum{5}-A]{Chun&Fessler:18TIP} that uses commutativity of convolution and Parseval's relation.\nGiven the current estimates of $\\{ z_{l,k} : l=1,\\ldots,L, k=1,\\ldots,K \\}$, the filter update problem of \\R{sys:CAOL} is equivalent to\n\\be{\n\\label{sys:filter}\n\\argmin_{\\{d_k\\}} \\frac{1}{2} \\sum_{k=1}^K \\sum_{l=1}^L \\left\\| \\Psi_l d_k - z_{l,k} \\right\\|_2^2 + \\beta g (D),\n}\nwhere $D$ is defined in \\R{eq:D}, $\\Psi_l \\in \\bbC^{N \\times R}$ is defined by\n\\be{\n\\label{eq:Psi_l}\n\\Psi_l := \\left[ \\begin{array}{ccc} P_{B_1} \\hat{x}_l & \\ldots & P_{B_R} \\hat{x}_l \\end{array} \\right],\n}\n$P_{B_r} \\in \\bbC^{N \\times \\hat{N}}$ is the $r\\rth$ (rectangular) selection matrix that selects $N$ rows corresponding to the indices $B_r = \\{ r, \\ldots, r+N-1 \\}$ from $I_{\\hat{N}}$, $\\{ \\hat{x}_l \\in \\bbC^{\\hat{N}} : l=1,\\ldots,L \\}$ is a set of padded training data, $\\hat{N} = N+R-1$.\nNote that applying $\\Psi_l^H$ in \\R{eq:Psi_l} to a vector of size $N$ is analogous to calculating cross-correlation between $\\hat{x}_l$ and the vector, i.e., $( \\Psi_l^H \\hat{z}_{l,k} )_{r} = \\sum_{n=1}^N \\hat{x}_{n+r-1}^{*} (\\hat{z}_{l,k})_n$, $r=1,\\ldots,R$.\nIn general, $\\hat{(\\cdot)}$ denotes a padded signal vector.\n\n\n\n\n\n\n\n\n\n\\subsubsection{Majorizer Design} \\label{sec:filt:maj}\n\n\n\nThis subsection designs multiple majorizers for the $D$-block optimization and compares their required computational complexity and tightness. \nThe next proposition considers the structure of $\\Psi_l$ in \\R{eq:Psi_l} to obtain the Hessian $\\sum_{l=1}^L \\Psi^H_l \\Psi_l \\in \\bbC^{R \\times R}$ in \\R{sys:filter} for an arbitrary boundary condition.\n\\prop{[Exact Hessian matrix $M_D$] \\label{p:MDexhess}\nThe following matrix $M_D \\in \\bbC^{R \\times R}$ is identical to $\\sum_{l=1}^L \\Psi^H_l \\Psi_l$:\n\\be{\n\\label{eq:soln:filterHess}\n\\left[ M_D \\right]_{r,r'} = \\sum_{l=1}^L \\ip{P_{B_r} \\hat{x}_l}{P_{B_{r'}} \\hat{x}_l}, \\quad r,r' = 1,\\ldots,R.\n}\n}\n\n\nA sufficiently large number of training signals (with $N \\geq R$), $L$, can guarantee $M_D = \\sum_{l=1}^L \\Psi^H_l \\Psi_l \\succ 0$ in Proposition~\\ref{p:MDexhess}.\nThe drawback of using Proposition~\\ref{p:MDexhess} is its polynomial computational complexity, i.e., $O (L R^2 N)$ -- see Table \\ref{tab:MD:compt}. \nWhen $L$ (the number of training signals) or $N$ (the size of training signals) are large, the quadratic complexity with the size of filters -- $R^2$ -- can quickly increase the total computational costs when multiplied by $L$ and $N$. (The BPG setup in \\cite{Xu&Yin:17JSC} additionally requires $O (R^3)$ because it uses the eigendecomposition of \\R{eq:soln:filterHess} to calculate the Lipschitz constant.)\n\n\n\n\n\n\nConsidering CAOL problems \\R{sys:CAOL} themselves, different from CDL \\cite{Chun&Fessler:18TIP, Chun&Fessler:17SAMPTA, Wohlberg:16TIP, Heide&eta:15CVPR, Bristow&etal:13CVPR}, the complexity $O (L R^2 N)$ in applying Proposition~\\ref{p:MDexhess} is reasonable. In BPEG-M-based CDL \\cite{Chun&Fessler:18TIP, Chun&Fessler:17SAMPTA}, a majorization matrix for kernel update is calculated every iteration because it depends on updated sparse codes; however, in CAOL, one can precompute $M_D$ via Proposition~\\ref{p:MDexhess} (or Lemmas~\\ref{l:MDdiag}--\\ref{l:MDscaleI} below) without needing to change it every kernel update.\nThe polynomial computational cost in applying Proposition~\\ref{p:MDexhess} becomes problematic only when the training signals change.\nExamples include \\textit{1)} hierarchical CAOL, e.g., CNN in Appendix~\\ref{sec:CNN}, \\textit{2)} \\dquotes{adaptive-filter MBIR} particularly with high-dimensional signals \\cite{Cai&etal:14ACHA, Xu&etal:12TMI, Elad&Aharon:06TIP}, and \\textit{3)} online learning \\cite{Liu&etal:17arXiv, Mairal&etal:09ICML}. Therefore, we also describe a more efficiently computable majorization matrix at the cost of looser bounds (i.e., slower convergence; see Fig~\\ref{fig:Comp:diffBPGM}). \nApplying Lemma~\\ref{l:diag(|At|W|A|1)}, we first introduce a diagonal majorization matrix $M_D$ for the Hessian $\\sum_{l} \\Psi_l^H \\Psi_l$ in \\R{sys:filter}:\n\n\n\n\\begin{table}[!pt]\t\n\n\\centering\n\\renewcommand{\\arraystretch}{1.1}\n\t\n\\caption{Computational complexity of different majorization matrix designs for the filter update problem \\R{sys:filter}}\t\n\\label{tab:MD:compt}\n\t\n\\begin{tabular}{C{2.2cm}C{2cm}}\n\\hline \\hline\nLemmas~\\ref{l:MDdiag}--\\ref{l:MDscaleI} & Proposition~\\ref{p:MDexhess} \\\\ \n\\hline\n$O( L R N )$ & $O( L R^2 N )$ \\\\\n\\hline \\hline\n\\end{tabular}\n\\end{table}\n\n\n\n\\lem{[Diagonal majorization matrix $M_D$] \\label{l:MDdiag}\nThe following matrix $M_D \\in \\bbC^{R \\times R}$ satisfies $M_D \\succeq \\sum_{l=1}^L \\Psi^H_l \\Psi_l$:\n\\be{\n\\label{eq:soln:filterMaj:loose}\nM_D = \\diag \\! \\left( \\sum_{l=1}^L | \\Psi_l^H | | \\Psi_l | 1_{R} \\right),\n}\nwhere $|\\cdot|$ takes the absolute values of the elements of a matrix. \n}\nThe majorization matrix design in Lemma~\\ref{l:MDdiag} is more efficient to compute than that in Proposition~\\ref{p:MDexhess}, because no $R^2$-factor is needed for calculating $M_D$ in Lemma~\\ref{l:MDdiag}, i.e., $O( L R N )$; see Table~\\ref{tab:MD:compt}. \nDesigning $M_D$ in Lemma~\\ref{l:MDdiag} takes fewer calculations than \\cite[Lem.~5.1]{Chun&Fessler:18TIP} using Fourier approaches, when $R \\!<\\! \\log ( \\hat{N} )$.\nUsing Lemma~\\ref{l:diag(|A|1)}, we next design a potentially sharper majorization matrix than \\R{eq:soln:filterMaj:loose}, while maintaining the cost $O( L R N )$:\n\\lem{[Scaled identity majorization matrix $M_D$] \\label{l:MDscaleI}\nThe following matrix $M_D \\in \\bbC^{R \\times R}$ satisfies $M_D \\succsim \\sum_{l=1}^L \\Psi^H_l \\Psi_l$:\n\\be{\n\\label{eq:soln:filterMaj}\nM_D = \\sum_{r=1}^R \\left| \\sum_{l=1}^L \\ip{P_{B_1} \\hat{x}_l}{P_{B_{r}} \\hat{x}_l} \\right| \\cdot I_R,\n}\nfor a circular boundary condition.\n}\n\\prf{\\renewcommand{\\qedsymbol}{}\nSee Section~\\ref{sec:prf:l:MD} of the supplementary material.\n}\n\n\n\n\\begin{figure}[!t]\n\\centering\n\\begin{tabular}{c}\n\\vspace{-0.25em}\\includegraphics[scale=0.5, trim=0 0.2em 2.2em 1.2em, clip]{.\/Fig_v1\/obj_detDvsrandD_fruit_v1.eps} \\\\\n{\\small (a) The fruit dataset ($L = 10$, $N = 100 \\!\\times\\! 100$)} \\\\\n\\vspace{-0.25em}\\includegraphics[scale=0.5, trim=0 0.2em 2.2em 1em, clip]{.\/Fig_v1\/obj_detDvsrandD_city_v1.eps} \\\\\n{\\small (b) The city dataset ($L = 10$, $N = 100 \\!\\times\\! 100$)} \\\\\n\\end{tabular}\n\n\\vspace{-0.25em}\n\\caption{Cost minimization comparisons in CAOL \\R{sys:CAOL:orth} with different BPG-type algorithms and datasets ($R \\!=\\! K \\!=\\! 49$ and $\\alpha \\!=\\! 2.5 \\!\\times\\! 10^{-4}$; solution \\R{eq:soln:spCode:exact} was used for sparse code updates; BPG (Xu \\& Ying '17) \\cite{Xu&Yin:17JSC} used the maximum eigenvalue of Hessians for Lipschitz constants; the cross mark \\text{\\sffamily x} denotes a termination point).\nA sharper majorization leads to faster convergence of BPEG-M;\nfor all the training datasets considered in this paper, the majorization matrix in Proposition~\\ref{p:MDexhess} is sharper than those in Lemmas~\\ref{l:MDdiag}--\\ref{l:MDscaleI}. \n}\n\\label{fig:Comp:diffBPGM}\n\\end{figure}\n\n\n\n\\begin{figure}[!tp]\n\\centering\n\\begin{tabular}{c}\n\\vspace{-0.25em}\\includegraphics[scale=0.58, trim=0 0.2em 2em 1em, clip]{.\/Fig_v1\/obj_diffBPGM_fruit_paper_v2.eps} \\\\\n{\\small (a) The fruit dataset ($L = 10$, $N = 100 \\!\\times\\! 100$)} \\\\\n\\vspace{-0.25em}\\includegraphics[scale=0.58, trim=0 0.2em 2em 1em, clip]{.\/Fig_v1\/obj_diffBPGM_city_paper_v2.eps} \\\\\n{\\small (b) The city dataset ($L = 10$, $N = 100 \\!\\times\\! 100$)} \\\\\n\\end{tabular}\n\n\\vspace{-0.25em}\n\\caption{Cost minimization comparisons in CAOL \\R{sys:CAOL:orth} with different BPEG-M algorithms and datasets (Lemma~\\ref{l:MDdiag} was used for $M_D$; $R \\!=\\! K \\!=\\! 49$; deterministic filter initialization and random sparse code initialization). \nUnder the sharp majorization regime, maintaining sharp majorization (i.e., $\\lambda_D \\!=\\! 1+\\epsilon$) provides faster convergence than giving more weight on extrapolation (i.e., $\\lambda_D \\!=\\! 2$).\n(The same behavior was found in sparse-view CT application \\cite[Fig.~3]{Chun&Fessler:18Asilomar}.)\nThere exist no differences in convergence between solution \\R{eq:soln:spCode:exact} and solution \\R{eq:soln:spCode} using $\\{ \\lambda_Z \\!=\\! 1+\\epsilon \\}$. \n}\n\\label{fig:Comp:exact_vs_approx}\n\\end{figure}\n\n\n\n\n\nFor all the training datasets used in this paper, we observed that the tightness of majorization matrices in Proposition~\\ref{p:MDexhess} and Lemmas~\\ref{l:MDdiag}--\\ref{l:MDscaleI} for the Hessian $\\sum_{l} \\Psi^H_l \\Psi_l$ is given by\n\\be{\n\\label{eq:MD:tight}\n\\sum_{l=1}^L \\Psi^H_l \\Psi_l = \\text{\\R{eq:soln:filterHess}} \\preceq \\text{\\R{eq:soln:filterMaj}} \\preceq \\text{\\R{eq:soln:filterMaj:loose}}.\n}\n(Note that \\R{eq:soln:filterHess}\\,$\\preceq$\\,\\R{eq:soln:filterMaj:loose} always holds regardless of training data.)\nFig.~\\ref{fig:Comp:diffBPGM} illustrates the effects of the majorizer sharpness in \\R{eq:MD:tight} on CAOL convergence rates.\nAs described in Section~\\ref{sec:BPGM:setup}, selecting $\\lambda_D$ (see \\R{sys:prox:filter:orth} and \\R{sys:prox:filter:div} below) controls the tradeoff between majorization sharpness and extrapolation effect. \nWe found that using fixed $\\lambda_D = 1+\\epsilon$ \ngives faster convergence than $\\lambda_D = 2 $; see Fig.~\\ref{fig:Comp:exact_vs_approx} (this behavior is more obvious in solving the CT MBIR model in \\R{sys:CT&CAOL} via BPEG-M -- see \\cite[Fig.~3]{Chun&Fessler:18Asilomar}).\nThe results in Fig.~\\ref{fig:Comp:exact_vs_approx} and \\cite[Fig.~3]{Chun&Fessler:18Asilomar} show that, under the sharp majorization regime, maintaining sharper majorization is more critical in accelerating the convergence of BPEG-M than giving more weight to extrapolation.\n\n\n\nSections \\ref{sec:prox:filter:orth} and \\ref{sec:prox:filter:div} below apply the majorization matrices designed in this section to proximal mappings of $D$-optimization in \\R{sys:CAOL:orth} and \\R{sys:CAOL:div}, respectively.\n\n\n\n\n\n\n\n\n\n\\subsubsection{Proximal Mapping with Orthogonality Constraint} \\label{sec:prox:filter:orth}\n\nThe corresponding proximal mapping problem of \\R{sys:filter} using the orthogonality constraint in \\R{sys:CAOL:orth} is given by\n\\ea{\n\\label{sys:prox:filter:orth}\n\\{ d_k^{(i+1)} \\} = \\argmin_{\\{ d_k \\}} &~ \n\\sum_{k=1}^K \\frac{1}{2} \\left\\| d_k - \\nu_k^{(i+1)} \\right\\|_{\\widetilde{M}_D}^2, \\nn\n\\\\\n\\mathrm{subject~to} &~ D D^H = \\frac{1}{R} \\cdot I,\n}\nwhere\n\\ea{\n\\nu_{k}^{(i+1)} &= \\textstyle \\acute{d}_k^{(i+1)} - \\widetilde{M}_D^{-1} \\sum_{l=1}^L \\Psi_l^H \\Big( \\Psi_l \\acute{d}_k^{(i+1)} - z_{l,k} \\Big),\n\\label{eq:nuk} \n\\\\\n\\acute{d}_k^{(i+1)} &= d_k^{(i)} + E_D^{(i+1)} \\Big( d_k^{(i)} - d_k^{(i-1)} \\Big), \\label{eq:dk_ac}\n}\nfor $k=1,\\ldots,K$, and $\\widetilde{M}_D = \\lambda_D M_D$ by \\R{update:Mtilde}. \nOne can parallelize over $k = 1,\\ldots,K$ in computing $\\{ \\nu_k^{(i+1)} \\}$ in \\R{eq:nuk}.\nThe proposition below provides an optimal solution to \\R{sys:prox:filter:orth}:\n\n\\prop{\n\\label{p:orth}\nConsider the following constrained minimization problem:\n\\begingroup\n\\setlength{\\thinmuskip}{1.5mu}\n\\setlength{\\medmuskip}{2mu plus 1mu minus 2mu}\n\\setlength{\\thickmuskip}{2.5mu plus 2.5mu}\n\\ea{\n\\label{p:eq:TF}\n\\min_{D} ~ \\nm{ \\widetilde{M}_D^{1\/2} D - \\widetilde{M}_D^{1\/2} \\mathcal{V} }_{\\mathrm{F}}^2, \\quad \n\\mathrm{subj.~to} ~ D D^H = \\frac{1}{R} \\cdot I,\n}\n\\endgroup\nwhere $D$ is given as \\R{eq:D}, $\\mathcal{V} = [\\nu_{1}^{(i+1)} \\cdots \\nu_{K}^{(i+1)}] \\in \\bbC^{R \\times K}$, $\\widetilde{M}_D = \\lambda_D M_D$, and $M_D \\in \\bbR^{R \\times R}$ is given by \\R{eq:soln:filterHess}, \\R{eq:soln:filterMaj:loose}, or \\R{eq:soln:filterMaj}.\nThe optimal solution to \\R{p:eq:TF} is given by\n\\bes{\n D^\\star = \\frac{1}{\\sqrt{R}} \\cdot U \\left[ \\arraycolsep=1pt \\begin{array}{cc} I_R, & 0_{R \\times (K-R)} \\end{array} \\right] V^H, \\quad \\mbox{for}~ R \\leq K,\n}\nwhere $\\widetilde{M}_D \\mathcal{V}$ has (full) singular value decomposition, $\\widetilde{M}_D \\mathcal{V} = U \\Lambda V^H$.\n}\n\\prf{\\renewcommand{\\qedsymbol}{}\nSee Section~\\ref{sec:prf:p:orth} of the supplementary material.\n}\n\n\nWhen using Proposition~\\ref{p:MDexhess}, $\\widetilde{M}_D \\nu_{k}^{(i+1)}$ of $\\widetilde{M}_D \\mathcal{V}$ in Proposition~\\ref{p:orth} simplifies to the following update:\n\\bes{\n\\widetilde{M}_D \\nu_{k}^{(i+1)} = (\\lambda_D - 1) M_D \\acute{d}_k^{(i+1)} + \\sum_{l=1}^L \\Psi_l^H z_{l,k}.\n}\nSimilar to obtaining $\\{ \\nu_k^{(i+1)} \\}$ in \\R{eq:nuk}, computing $\\{ \\widetilde{M}_D \\nu_k^{(i+1)} : k=1,\\ldots,K \\}$ is parallelizable over $k$.\n\n\n\n\n\n\n\n\n\\subsubsection{Proximal Mapping with Diversity Promoting Regularizer} \\label{sec:prox:filter:div}\n\nThe corresponding proximal mapping problem of \\R{sys:filter} using the norm constraint and diversity promoting regularizer in \\R{sys:CAOL:div} is given by\n\\ea{\n\\label{sys:prox:filter:div}\n\\{ d_k^{(i+1)} \\} = \\argmin_{\\{ d_k \\}} &~ \n\\sum_{k=1}^K \\frac{1}{2} \\left\\| d_k - \\nu_k^{(i+1)} \\right\\|_{\\widetilde{M}_D}^2 + \\frac{\\beta}{2} g_{\\text{div}}( D ), \\nn\n\\\\\n\\mathrm{subject~to} &~ \\left\\| d_k \\right\\|_2^2 = \\frac{1}{R}, \\quad k = 1,\\ldots, K,\n}\nwhere $g_{\\text{div}}( D )$, $\\nu_{k}^{(i+1)}$, and $\\acute{d}_k^{(i+1)}$ are given as in \\R{sys:CAOL:div}, \\R{eq:nuk}, and \\R{eq:dk_ac}, respectively.\nWe first decompose the regularization term $g_{\\text{div}}(D)$ as follows:\n\\begingroup\n\\allowdisplaybreaks\n\\ea{\n\\label{eq:DtD}\ng_{\\text{div}} (D)\n& = \\sum_{k=1}^K \\sum_{k'=1}^K \\big( d_k^H d_{k'} d_{k'}^H d_k - R^{-1} \\big) \\nn\n\\\\\n& = \\sum_{k=1}^K d_k^H \\bigg( \\sum_{k' \\neq k} d_{k'} d_{k'}^H \\bigg) d_k + \\big( d_k^H d_k - R^{-1} \\big)^2 \\nn\n\\\\\n& = \\sum_{k=1}^K d_k^H \\Gamma_k d_k, \n}\n\\endgroup\nwhere the equality in \\R{eq:DtD} holds by using the constraint in \\R{sys:prox:filter:div}, and the Hermitian matrix $\\Gamma_k \\in \\bbC^{R \\times R}$ is defined by\n\\be{\n\\label{eq:Ek}\n\\Gamma_k := \\sum_{k' \\neq k} d_{k'} d_{k'}^H.\n}\nUsing \\R{eq:DtD} and \\R{eq:Ek}, we rewrite \\R{sys:prox:filter:div} as\n\\ea{\n\\label{sys:prox:filter2}\nd_k^{(i+1)} = \\argmin_{d_k} &~ \n\\frac{1}{2} \\left\\| d_k - \\nu_k^{(i)} \\right\\|_{\\widetilde{M}_D}^2 + \\frac{\\beta}{2} d_k^H \\Gamma_k d_k, \\nn\n\\\\\n\\mathrm{subject~to} &~ \\left\\| d_k \\right\\|_2^2 = \\frac{1}{R}, \\quad k = 1,\\ldots,K.\n}\nThis is a quadratically constrained quadratic program with $\\{ \\widetilde{M}_D + \\beta \\Gamma_k \\succ 0 : k = 1,\\ldots,K\\}$. \nWe apply an accelerated Newton's method to solve \\R{sys:prox:filter2}; see Section~\\ref{sec:Newton}.\nSimilar to solving \\R{sys:prox:filter:orth} in Section~\\ref{sec:prox:filter:orth}, solving \\R{sys:prox:filter:div} is a small-dimensional problem ($K$ separate problems of size $R$). \n\n\n\n\n\\subsection{Sparse Code Update: $\\{z_{l,k} \\}$-Block Optimization} \\label{sec:CAOL:spCd}\n\nGiven the current estimate of $D$, the sparse code update problem for \\R{sys:CAOL} is given by\n\\be{\n\\label{sys:spCode}\n\\argmin_{\\{ z_{l,k} \\}} \\sum_{l=1}^L \\sum_{k=1}^K \\frac{1}{2} \\left\\| d_k \\circledast x_{l} - z_{l,k} \\right\\|_2^2 + \\alpha \\left\\| z_{l,k} \\right\\|_0.\n}\nThis problem separates readily, \nallowing parallel computation with $LK$ threads.\nAn optimal solution to \\R{sys:spCode} is efficiently obtained by the well-known hard thresholding:\n\\be{\n\\label{eq:soln:spCode:exact}\nz_{l,k}^{(i+1)} = \\cH_{\\!\\sqrt{2\\alpha}} \\left( d_k \\circledast x_{l} \\right), \n}\nfor $k=1,\\ldots,K$ and $l=1,\\ldots,L$, where \n\\be{\n\\label{eq:def:hardthr}\n\\cH_a (x)_{n} := \\left\\{ \\begin{array}{cc} 0, & | x_{n} | < a_{n}, \\\\ x_{n}, & | x_{n} | \\geq a_{n}. \\end{array} \\right.\n}\nfor all $n$. \nConsidering $\\lambda_Z$ (in $\\widetilde{M}_Z = \\lambda_Z M_Z$) as $\\lambda_Z \\!\\rightarrow\\! 1$, the solution obtained by the BPEG-M approach becomes equivalent to \\R{eq:soln:spCode:exact}. \nTo show this, observe first that the BPEG-M-based solution (using $M_Z = I_{N}$) to \\R{sys:spCode} is obtained by\n\\ea{\n\\label{eq:soln:spCode}\nz_{l,k}^{(i+1)} &= \\cH_{\\!\\sqrt{ \\frac{2\\alpha}{\\lambda_Z}}} \\Big( \\zeta_{l,k}^{(i+1)} \\Big),\n\\\\\n\\zeta_{l,k}^{(i+1)} &= \\left( 1 - \\lambda_Z^{-1} \\right) \\cdot \\acute{z}_{l,k}^{(i+1)} + \\lambda_Z^{-1} \\cdot d_k \\circledast x_{l},\n\\nn \\\\\n\\acute{z}_{l,k}^{(i+1)} &= z_{l,k}^{(i)} + E_Z^{(i+1)} \\Big( z_{l,k}^{(i)} - z_{l,k}^{(i-1)} \\Big). \\nn\n}\nThe downside of applying solution \\R{eq:soln:spCode} is that it would require additional memory to store the corresponding extrapolated points -- $\\{ \\acute{z}_{l,k}^{(i+1)} \\}$ -- and the memory grows with $N$, $L$, and $K$.\nConsidering the sharpness of the majorizer in \\R{sys:spCode}, i.e., $M_Z = I_{N}$, and the memory issue, it is reasonable to consider the solution \\R{eq:soln:spCode} with no extrapolation, i.e., $\\{ E_Z^{(i+1)} = 0 \\}$:\n\\bes{\nz_{l,k}^{(i+1)} \n= \\cH_{\\! \\sqrt{\\frac{2\\alpha}{\\lambda_Z}}} \\Big( (\\lambda_Z - 1)^{-1}\\lambda_Z \\cdot z_{l,k}^{(i)} + \\lambda_Z^{-1} \\cdot d_k \\circledast x_{l} \\Big)\n}\nbecoming equivalent to \\R{eq:soln:spCode:exact} as $\\lambda_Z \\!\\rightarrow\\! 1$.\n\n\n\n\n\n\nSolution \\R{eq:soln:spCode:exact} has two benefits over \\R{eq:soln:spCode}: compared to \\R{eq:soln:spCode}, \\R{eq:soln:spCode:exact} requires only half the memory to update all $z_{l,k}^{(i+1)}$ vectors and no additional computations related to $\\acute{z}_{l,k}^{(i+1)}$.\nWhile having these benefits, empirically \\R{eq:soln:spCode:exact} has equivalent convergence rates as \\R{eq:soln:spCode} using $\\{ \\lambda_Z \\!=\\! 1+\\epsilon \\}$; see Fig.~\\ref{fig:Comp:exact_vs_approx}. Throughout the paper, we solve the sparse coding problems (e.g., \\R{sys:spCode} and $\\{ z_k \\}$-block optimization in \\R{sys:CT&CAOL}) via optimal solutions in the form of \\R{eq:soln:spCode:exact}.\n\n\n\n\n\n\n\\subsection{Lower Memory Use than Patch-Domain Approaches} \\label{sec:CAOL:memory}\n\nThe convolution perspective in CAOL \\R{sys:CAOL} requires much less memory than conventional patch-domain approaches; \nthus, it is more suitable for learning filters from large datasets or applying the learned filters to high-dimensional MBIR problems.\nFirst, consider the training stage (e.g., \\R{sys:CAOL}).\nThe patch-domain approaches, e.g., \\cite{Cai&etal:14ACHA, Ravishankar&Bressler:15TSP, Aharon&Elad&Bruckstein:06TSP}, require about $R$ times more memory to store training signals. For example, 2D patches extracted by $\\sqrt{R} \\!\\times\\! \\sqrt{R}$-sized windows (with \\dquotes{stride} one and periodic boundaries \\cite{Cai&etal:14ACHA, Papyan&Romano&Elad:17JMLR}, as used in convolution) require about $R$ (e.g., $R = 64$ \\cite{Aharon&Elad&Bruckstein:06TSP, Ravishankar&Bressler:15TSP}) times more memory than storing the original image of size $\\sqrt{N} \\!\\times\\! \\sqrt{N}$. For $L$ training images, their memory usage dramatically increases with a factor $L R N$.\nThis becomes even more problematic in forming hierarchical representations, e.g., CNNs -- see Appendix~\\ref{sec:CNN}.\nUnlike the patch-domain approaches, the memory use of CAOL \\R{sys:CAOL} only depends on the $L N$-factor to store training signals.\nAs a result, the BPEG-M algorithm for CAOL \\R{sys:CAOL:orth} requires about two times less memory than the patch-domain approach \\cite{Cai&etal:14ACHA} (using BPEG-M).\nSee Table~\\ref{tab:AOL}-B.\n(Both the corresponding BPEG-M algorithms use identical computations per iteration that scale with $L R^2 N$; see Table~\\ref{tab:AOL}-A.)\n\n\n\n\n\nSecond, consider solving MBIR problems. Different from the training stage, the memory burden depends on how one applies the learned filters.\nIn \\cite{Pfister&Bresler:17ICASSP}, the learned filters are applied with the conventional convolutional operators -- e.g., $\\circledast$ in \\R{sys:CAOL} -- and, thus, there exists no additional memory burden. \nHowever, in \\cite{Elad&Aharon:06TIP, Chun&etal:17Fully3D, Zheng&etal:19TCI}, the $\\sqrt{R} \\!\\times\\! \\sqrt{R}$-sized learned kernels are applied with a matrix constructed by many overlapping patches extracted from the updated image at each iteration. \nIn adaptive-filter MBIR problems \\cite{Elad&Aharon:06TIP, Cai&etal:14ACHA, Pfister&Bresler:15SPIE}, the memory issue pervades the patch-domain approaches.\n\n\n\n\n\\begin{table}[!t]\t\n\n\\centering\n\\renewcommand{\\arraystretch}{1.1}\n\t\n\\caption{\nComparisons of computational complexity and memory usages between CAOL and patch-domain approach\n}\t\n\\label{tab:AOL}\n\t\n\\begin{tabular}{C{2.15cm}|C{2.85cm}|C{2.3cm}}\n\\hline \\hline\n\\multicolumn{3}{c}{A.~Computational complexity per BPEG-M iteration}\n\\\\\n\\hline\n& Filter update & Sparse code update\n\\\\\n\\hline\nCAOL \\R{sys:CAOL:orth} & $O(L K R N) + O(R^2 K)$ & $O(L K R N)$ \n\\\\ \n\\hline\nPatch-domain \\cite{Cai&etal:14ACHA}$^\\dagger$ & $O (L R^2 N) + O (R^3)$ & $O (L R^2 N)$\n\\\\\n\\hline \\hline\n\\end{tabular}\n\n\\vspace{0.75pc}\n\n\\begin{tabular}{C{2.15cm}|C{2.3cm}|C{2.3cm}}\n\\hline \\hline\n\\multicolumn{3}{c}{B.~Memory usage for BPEG-M algorithm}\n\\\\\n\\hline\n& Filter update & Sparse code update\n\\\\\n\\hline\nCAOL \\R{sys:CAOL:orth} & $O(LN) + O(RK)$ & $O(L K N)$\n\\\\ \n\\hline\nPatch-domain \\cite{Cai&etal:14ACHA}$^\\dagger$ & $O(L R N) + O(R^2)$ & $O(L R N)$\n\\\\\n\\hline \\hline\n\\end{tabular}\n\n\\medskip\n\\begin{myquote}{0.1in}\n$^\\dagger$\nThe patch-domain approach \\cite{Cai&etal:14ACHA} considers the orthogonality constraint in \\R{sys:CAOL:orth} with $R \\!=\\! K$; see Section~\\ref{sec:CAOL:orth}.\nThe estimates consider all the extracted overlapping patches of size $R$ with the stride parameter $1$ and periodic boundaries, as used in convolution.\n\\end{myquote}\n\\end{table}\n\n\n\n\n\n\n\n\\section{Sparse-View CT MBIR using Convolutional Regularizer Learned via CAOL, and BPEG-M} \\label{sec:CTrecon}\n\nThis section introduces a specific example of applying the learned convolutional regularizer, i.e., $F(D^\\star, \\{ z_{l,k} \\})$ in \\R{sys:CAOL}, from a representative dataset to recover images in \\textit{extreme} imaging that collects highly undersampled or noisy measurements.\nWe choose a sparse-view CT application since it has interesting challenges in reconstructing images that include Poisson noise in measurements, nonuniform noise or resolution properties in reconstructed images, and complicated (or no) structures in the system matrices. \nFor CT, undersampling schemes can significantly reduce the radiation dose and cancer risk from CT scanning. \nThe proposed approach can be applied to other applications (by replacing the data fidelity and spatial strength regularization terms in \\R{sys:CT&CAOL} below).\n\n\n\nWe pre-learn TF filters $\\{ d_k^\\star \\in \\bbR^K : k = 1,\\ldots,K \\}$ via CAOL \\R{sys:CAOL:orth} with a set of high-quality (e.g., normal-dose) CT images $\\{ x_l : l = 1,\\ldots,L \\}$.\nTo reconstruct a linear attenuation coefficient image $x \\in \\bbR^{N'}$ from post-log measurement $y \\in \\bbR^m$ \\cite{Chun&Talavage:13Fully3D, Chun&etal:17Fully3D}, we apply the learned convolutional regularizer to CT MBIR and solve the following block multi-nonconvex problem \\cite{Chun&Fessler:18Asilomar, Crockett&etal:19CAMSAP}:\n\\begingroup\n\\allowdisplaybreaks\n\\setlength{\\thinmuskip}{1.5mu}\n\\setlength{\\medmuskip}{2mu plus 1mu minus 2mu}\n\\setlength{\\thickmuskip}{2.5mu plus 2.5mu}\n\\ea{\n\\label{sys:CT&CAOL}\n\\argmin_{x \\geq 0} & \\underbrace{ \\frac{1}{2} \\left\\| y - A x \\right\\|_W^2 }_{\\text{data fidelity $f(x;y)$}} + \n\\nn \\\\\n\\gamma \\cdot & \\underbrace{ \\min_{\\{ z_k \\}} \\sum_{k=1}^K \\frac{1}{2} \\left\\| d_{k}^\\star \\circledast x - z_k \\right\\|_2^2 \n+ \\alpha' \\sum_{n=1}^{N'} \\psi_{j} \\phi( (z_{k})_n ) }_{\\text{learned convolutional regularizer $g(x,\\{z_k\\}; \\{ d_k \\})$}}. \\tag{P3}\n}\n\\endgroup\n\\!\\!Here, $A \\in \\bbR^{m \\times N'}$ is a CT system matrix, $W \\in \\bbR^{m \\times m}$ is a (diagonal) weighting matrix with elements $\\{ W_{l,l} = \\rho_l^2 \/ ( \\rho_l + \\sigma^2 ) : l = 1,\\ldots,m \\}$ based on a Poisson-Gaussian model for the pre-log measurements $\\rho \\in \\bbR^m$ with electronic readout noise variance $\\sigma^2$ \\cite{Chun&Talavage:13Fully3D, Chun&etal:17Fully3D, Zheng&etal:19TCI}, $\\psi \\in \\bbR^{N'}$ is a pre-tuned spatial strength regularization vector \\cite{Fessler&Rogers:96TIP} with non-negative elements\n$\\{ \\psi_{n} = ( \\sum_{l=1}^m A_{l,n}^2 W_{l,l} )^{1\/2} \/ ( \\sum_{l=1}^m A_{l,n}^2 )^{1\/2} : n = 1,\\ldots,N' \\}$\\footnote{\nSee details of computing $\\{ A_{l,j}^2 : \\forall l,j \\}$ in \\cite{Chun&Fessler:18Asilomar}.\n} \nthat promotes uniform resolution or noise properties in the reconstructed image \\cite[Appx.]{Chun&etal:17Fully3D}, an indicator function $\\phi(a)$ is equal to $0$ if $a = 0$, and is $1$ otherwise, $z_k \\in \\bbR^{N'}$ is unknown sparse code for the $k\\text{th}$ filter, \nand $\\alpha' \\!>\\! 0$ is a thresholding parameter. \n\n\n\n\n\nWe solved \\R{sys:CT&CAOL} via reBPEG-M in Section~\\ref{sec:reBPG-M} with a two-block scheme \\cite{Chun&Fessler:18Asilomar}, and summarize the corresponding BPEG-M updates as\n\\begingroup\n\\allowdisplaybreaks\n\\ea{\n\\label{eq:soln:CT:bpgm}\nx^{(i+1)} &= \\bigg[ \\big( \\widetilde{M}_A + \\gamma I_R \\big)^{-1} \\cdot \\bigg( \\widetilde{M}_A \\eta^{(i+1)} + \n\\nn \\\\\n& \\hspace{1.95em} \\gamma \\sum_{k=1}^K ( P_f d_k^\\star ) \\circledast \\cH_{\\!\\sqrt{2 \\alpha' \\psi}} \\big( d_k^\\star \\circledast x^{(i)} \\big) \\bigg) \\bigg]_{\\geq 0},\n}\nwhere \n\\ea{\n\\eta^{(i+1)} &= \\acute{x}^{(i+1)} - \\widetilde{M}_A^{-1} A^T W \\Big( A \\acute{x}^{(i+1)} -y \\Big),\n\\label{eq:soln:CT:bpgm:eta} \\\\\n\\acute{x}^{(i+1)} &= x^{(i)} + E_A^{(i+1)} \\Big( x^{(i)} - x^{(i-1)} \\Big),\n\\nn\n}\n\\endgroup\n$\\widetilde{M}_A = \\lambda_A M_A$ by \\R{update:Mtilde}, a diagonal majorization matrix $M_A \\succeq A^T W A$ is designed by Lemma~\\ref{l:diag(|At|W|A|1)}, \nand $P_f \\in \\bbC^{R \\times R}$ flips a column vector in the vertical direction (e.g., it rotates 2D filters by $180^\\circ$).\nInterpreting the update \\R{eq:soln:CT:bpgm} leads to the following two remarks:\n\n\n\n\n\n\n\n\\rem{\n\\label{r:autoenc}\nWhen the convolutional regularizer learned via CAOL \\R{sys:CAOL:orth} is applied to MBIR, it works as an autoencoding CNN: \n\\be{\n\\label{eq:autoenc}\n\\cM ( x ) = \\sum_{k=1}^K (P_f d_k^\\star) \\circledast \\cH_{\\!\\sqrt{2 \\alpha_k'}} \\left( d_k^\\star \\circledast x \\right)\n}\n(setting $\\psi = 1_{N'}$ and generalizing $\\alpha'$ to $\\{ \\alpha'_k : k=1,\\ldots,K \\}$ in \\R{sys:CT&CAOL}).\nThis is an explicit mathematical motivation\nfor constructing architectures of iterative regression CNNs for MBIR, \ne.g., BCD-Net \\cite{Chun&Fessler:18IVMSP, Chun&etal:19MICCAI, Lim&etal:19TMI, Lim&etal:18NSSMIC}\nand Momentum-Net \\cite{Chun&etal:18Allerton, Chun&etal:18arXiv:momnet}. \nParticularly when the learned filters $\\{ d_k^\\star \\}$ in \\R{eq:autoenc} satisfy the TF condition, they are useful for compacting energy of an input signal $x$ and removing unwanted features via the non-linear thresholding in \\R{eq:autoenc}.\n}\n\n\n\n\\rem{\n\\label{r:autoenc-bpgm}\nUpdate \\R{eq:soln:CT:bpgm} improves the solution $x^{(i+1)}$ by weighting between \\textit{a)} the extrapolated point considering the data fidelity, i.e., $\\eta^{(i+1)}$ in \\R{eq:soln:CT:bpgm:eta}, and \\textit{b)} the \\dquotes{refined} update via the ($\\psi$-weighting) convolutional autoencoder, i.e., $\\sum_{k} ( P_f d_k^\\star ) \\circledast \\cH_{\\!\\sqrt{2 \\alpha' \\psi}} ( d_k^\\star \\circledast x^{(i)} )$.\n}\n\n\n\n\n\n\n\n\n\n\\section{Results and Discussion} \\label{sec:result}\n\n\\subsection{Experimental Setup} \\label{sec:exp}\n\nThis section examines the performance (e.g., scalability, convergence, and acceleration) and behaviors (e.g., effects of model parameters on filters structures and effects of dimensions of learned filter on MBIR performance) of the proposed CAOL algorithms and models, respectively.\n\n\\subsubsection{CAOL} \\label{sec:exp:CAOL}\n\nWe tested the introduced CAOL models\/algorithms for four datasets: \n\\textit{1)} the fruit dataset with $L = 10$ and $N = 100 \\!\\times\\! 100$ \\cite{Zeiler&etal:10CVPR}; \n\\textit{2)} the city dataset with $L = 10$ and $N = 100 \\!\\times\\! 100$ \\cite{Heide&eta:15CVPR}; \n\\textit{3)} the CT dataset of $L = 80$ and $N = 128 \\!\\times\\! 128$, created by dividing down-sampled $512 \\!\\times\\! 512$ XCAT phantom slices \\cite{Segars&etal:08MP} into $16$ sub-images \\cite{Olshausen&Field:96Nature, Bristow&etal:13CVPR} -- referred to the CT-(\\romnum{1}) dataset; \n\\textit{4)} the CT dataset of with $L = 10$ and $N = 512 \\!\\times\\! 512$ from down-sampled $512 \\!\\times\\! 512$ XCAT phantom slices \\cite{Segars&etal:08MP} -- referred to the CT-(\\romnum{2}) dataset.\nThe preprocessing includes intensity rescaling to $[0,1]$ \\cite{Zeiler&etal:10CVPR, Bristow&etal:13CVPR, Heide&eta:15CVPR} and\/or (global) mean substraction \\cite[\\S 2]{Jarrett&etal:09ICCV}, \\cite{Aharon&Elad&Bruckstein:06TSP}, as conventionally used in many sparse coding studies, e.g., \\cite{Aharon&Elad&Bruckstein:06TSP, Jarrett&etal:09ICCV, Zeiler&etal:10CVPR, Bristow&etal:13CVPR, Heide&eta:15CVPR}. \nFor the fruit and city datasets, we trained $K = 49$ filters of size $R = 7 \\!\\times\\! 7$.\nFor the CT dataset~(\\romnum{1}), we trained filters of size $R = 5 \\!\\times\\! 5$, with $K = 25$ or $K = 20$.\nFor CT reconstruction experiments, we learned the filters from the CT-(\\romnum{2}) dataset; however, we did not apply mean subtraction because it is not modeled in \\R{sys:CT&CAOL}. \n\nThe parameters for the BPEG-M algorithms were defined as follows.\\footnote{The remaining BPEG-M parameters not described here are identical to those in \\cite[\\Romnum{7}-A2]{Chun&Fessler:18TIP}.}\nWe set the regularization parameters $\\alpha, \\beta$ as follows:\n\\begin{itemize}\n\\item {CAOL \\R{sys:CAOL:orth}}: To investigate the effects of $\\alpha$, we tested \\R{sys:CAOL:orth} with different $\\alpha$'s in the case $R = K$. For the fruit and city datasets, we used $\\alpha = 2.5 \\!\\times\\! \\{ 10^{-5}, 10^{-4} \\}$; for the CT-(\\romnum{1}) dataset, we used $\\alpha = \\{ 10^{-4}, 2 \\!\\times\\! 10^{-3} \\}$. \nFor the CT-(\\romnum{2}) dataset (for CT reconstruction experiments), see details in \\cite[Sec.~\\Romnum{5}1]{Chun&Fessler:18Asilomar}.\n\n\\item {CAOL \\R{sys:CAOL:div}}: Once $\\alpha$ is fixed from the CAOL \\R{sys:CAOL:orth} experiments above, we tested \\R{sys:CAOL:div} with different $\\beta$'s to see its effects in the case $R > K$. For the CT-(\\romnum{1}) dataset, we fixed $\\alpha = 10^{-4}$, and used $\\beta = \\{ 5 \\!\\times\\! 10^6, 5 \\!\\times\\! 10^4 \\}$.\n\\end{itemize}\nWe set $\\lambda_D = 1+\\epsilon$ as the default.\nWe initialized filters in either deterministic or random ways.\nThe deterministic filter initialization follows that in \\cite[Sec.~3.4]{Cai&etal:14ACHA}.\nWhen filters were randomly initialized, we used a scaled one-vector for the first filter.\nWe initialize sparse codes mainly with a deterministic way that applies \\R{eq:soln:spCode:exact} based on $\\{ d_k^{(0)} \\}$.\nIf not specified, we used the random filter and deterministic sparse code initializations.\nFor BPG \\cite{Xu&Yin:17JSC}, we used the maximum eigenvalue of Hessians for Lipschitz constants in \\R{sys:filter}, and applied the gradient-based restarting scheme in Section~\\ref{sec:reBPGM}.\nWe terminated the iterations if the relative error stopping criterion (e.g., \\cite[(44)]{Chun&Fessler:18TIP}) is met before reaching the maximum number of iterations.\nWe set the tolerance value as $10^{-13}$ for the CAOL algorithms using Proposition~\\ref{p:MDexhess}, and $10^{-5}$ for those using Lemmas~\\ref{l:MDdiag}--\\ref{l:MDscaleI}, and the maximum number of iterations to $2 \\!\\times\\! 10^4$.\n\nThe CAOL experiments used the convolutional operator learning toolbox~\\cite{chun:19:convolt}.\n\n\n\n\n\\subsubsection{Sparse-View CT MBIR with Learned Convolutional Regularizer via CAOL} \\label{sec:exp:CT}\n\nWe simulated sparse-view sinograms of size $888 \\times 123$ (\\quotes{detectors or rays} $\\times$ \\quotes{regularly spaced projection views or angles}, where $984$ is the number of full views) with GE LightSpeed fan-beam geometry corresponding to a monoenergetic source with $10^5$ incident photons per ray and no background events, and electronic noise variance $\\sigma^2 \\!=\\! 5^2$. \nWe avoided an inverse crime in our imaging simulation and reconstructed images with a coarser grid with $\\Delta_x \\!=\\! \\Delta_y \\!=\\! 0.9766$ mm; see details in \\cite[Sec.~\\Romnum{5}-A2]{Chun&Fessler:18Asilomar}.\n\n\n\n\nFor EP MBIR, we finely tuned its regularization parameter to achieve both good root mean square error (RMSE) and structural similarity index measurement \\cite{Wang&etal:04TIP} values. \nFor the CT MBIR model \\R{sys:CT&CAOL}, we chose the model parameters $\\{ \\gamma, \\alpha' \\}$ that showed a good tradeoff between the data fidelity term and the learned convolutional regularizer, and set $\\lambda_A \\!=\\! 1+\\epsilon$.\nWe evaluated the reconstruction quality by the RMSE (in a modified Hounsfield unit, HU, where air is $0$ HU and water is $1000$ HU) in a region of interest.\nSee further details in \\cite[Sec.~\\Romnum{5}-A2]{Chun&Fessler:18Asilomar} and Fig.~\\ref{fig:CTrecon}.\n\n\nThe imaging simulation and reconstruction experiments used the Michigan image reconstruction toolbox~\\cite{fessler:16:irt}.\n\n\n\n\n\n\n\n\n\\subsection{CAOL with BPEG-M} \\label{sec:result:CAOL}\n\n\nUnder the sharp majorization regime (i.e., partial or all blocks have sufficiently tight bounds in Lemma~\\ref{l:QM}), the proposed convergence-guaranteed BPEG-M can achieve significantly faster CAOL convergence rates compared with the state-of-the-art BPG algorithm \\cite{Xu&Yin:17JSC} for solving block multi-nonconvex problems, by several generalizations of BPG (see Remark~\\ref{r:BPGM}) and two majorization designs (see Proposition~\\ref{p:MDexhess} and Lemma~\\ref{l:MDscaleI}). See Fig.~\\ref{fig:Comp:diffBPGM}.\nIn controlling the tradeoff between majorization sharpness and extrapolation effect of BPEG-M (i.e., choosing $\\{ \\lambda_b \\}$ in \\R{update:Mtilde}--\\R{update:Eb}), maintaining majorization sharpness is more critical than gaining stronger extrapolation effects to accelerate convergence under the sharp majorization regime.\nSee Fig.~\\ref{fig:Comp:exact_vs_approx}. \n\n\nWhile using about two times less memory (see Table~\\ref{tab:AOL}), CAOL \\R{sys:CAOL} learns TF filters corresponding to those given by the patch-domain TF learning in \\cite[Fig.~2]{Cai&etal:14ACHA}.\nSee Section~\\ref{sec:CAOL:memory} and Fig.~\\ref{fig:filters_BPGM} with deterministic $\\{ d_k^{(0)} \\}$.\nNote that BPEG-M-based CAOL \\R{sys:CAOL} requires even less memory than BPEG-M-based CDL in \\cite{Chun&Fessler:18TIP}, by using exact sparse coding solutions (e.g., \\R{eq:soln:spCode:exact} and \\R{eq:soln:CT:bpgm}) without saving their extrapolated points.\nIn particular, \nwhen tested with the large CT dataset of $\\{ L \\!=\\! 40, N \\!=\\! 512 \\!\\times\\! 512 \\}$,\nthe BPEG-M-based CAOL algorithm ran fine,\nwhile BPEG-M-based CDL \\cite{Chun&Fessler:18TIP} and patch-domain AOL \\cite{Cai&etal:14ACHA} were terminated due to exceeding available memory.\\footnote{\nTheir double-precision MATLAB implementations were tested on 3.3 GHz Intel Core i5 CPU with 32 GB RAM.\n}\nIn addition, the CAOL models \\R{sys:CAOL:orth} and \\R{sys:CAOL:div} are easily parallelizable with $K$ threads.\nCombining these results, the BPEG-M-based CAOL is a reasonable choice for learning filters from large training datasets. \nFinally, \\cite{Chun&etal:19SPL} shows theoretically how using many samples can improve CAOL, \naccentuating the benefits of the low memory usage of CAOL.\n\n\n\nThe effects of parameters for the CAOL models are shown as follows.\nIn CAOL~\\R{sys:CAOL:orth}, as the thresholding parameter $\\alpha$ increases, the learned filters have more elongated structures; see Figs.~\\ref{fig:CAOL:filters:CT}(a) and \\ref{fig:CAOL:filters:fruit&city&ct-ii}.\nIn CAOL~\\R{sys:CAOL:div}, when $\\alpha$ is fixed, increasing the filter diversity promoting regularizer $\\beta$ successfully lowers coherences between filters (e.g., $g_{\\text{div}}(D)$ in \\R{sys:CAOL:div}); see Fig.~\\ref{fig:CAOL:filters:CT}(b).\n\n\n\n\\begin{figure}[!t]\n\\small\\addtolength{\\tabcolsep}{-4pt}\n\\centering\n\n\\begin{tabular}{cc}\n\n\\vspace{-0.1em}\\hspace{-0.05em}\\includegraphics[scale=0.5, trim=2em 5.2em 0em 4em, clip]{.\/Fig\/filt_tfMexhess_CTeg1.eps} &\n\\vspace{-0.1em}\\hspace{-0.4em}\\includegraphics[scale=0.5, trim=2em 5.2em 0.1em 4em, clip]{.\/Fig\/filt_tfMexhess_CTeg2.eps} \\\\\n\n{\\small (a1) $\\alpha = 10^{-4}$} & {\\small (a2) $\\alpha = 2 \\!\\times\\! 10^{-3}$} \\\\\n\n\\multicolumn{2}{c}{\\small (a) Learned filters via CAOL \\R{sys:CAOL:orth} ($R = K = 25$)} \\\\\n\n\\includegraphics[scale=0.506, trim=2em 1.0em 0em 0.4em, clip]{.\/Fig\/filt_divMexhess_CTeg2-1.eps} &\n\\includegraphics[scale=0.506, trim=2em 1.0em 0em 0.4em, clip]{.\/Fig\/filt_divMexhess_CTeg1-1.eps} \\vspace{-0.3em} \\\\\n\n\\small{$g_{\\text{div}}(D) = \\mathbf{8.96 \\!\\times\\! 10^{-6}}$} & \\small{$g_{\\text{div}}(D) = \\mathbf{0.12}$} \\\\\n\n\\small{(b1) $\\alpha = 10^{-4}$, $\\beta = 5 \\!\\times\\! 10^{6}$} & {\\small (b2) $\\alpha = 10^{-4}$, $\\beta = 5 \\!\\times\\! 10^{4}$} \\\\\n\n\\multicolumn{2}{c}{\\small (b) Learned filters via CAOL \\R{sys:CAOL:div} ($R = 25, K = 20$) } \n\n\\end{tabular}\n\n\\vspace{-0.25em}\n\\caption{Examples of learned filters with different CAOL models and parameters (Proposition~\\ref{p:MDexhess} was used for $M_D$; the CT-(\\romnum{1}) dataset with a symmetric boundary condition).}\n\\label{fig:CAOL:filters:CT}\n\\end{figure}\n\n\n\n\nIn adaptive MBIR (e.g., \\cite{Elad&Aharon:06TIP, Cai&etal:14ACHA, Pfister&Bresler:15SPIE}), one may apply adaptive image denoising \\cite{Donoho:95TIT, Donoho&Johnstone:95JASA, Chang&Yu&Vetterli:00TIP, Blu&Luisier:07TIP, Liu&etal:15CVPR, Pfister&Bresler:17ICASSP} to optimize thresholding parameters. \nHowever, if CAOL \\R{sys:CAOL} and testing the learned convolutional regularizer to MBIR (e.g., \\R{sys:CT&CAOL}) are separated, selecting \\dquotes{optimal} thresholding parameters in (unsupervised) CAOL is challenging -- similar to existing dictionary or analysis operator learning methods.\nOur strategy to select the thresholding parameter $\\alpha$ in CAOL~\\R{sys:CAOL:orth} (with $R=K$) is given as follows.\nWe first apply the first-order finite difference filters $\\{ d_k : \\| d_k \\|_2^2 = 1\/R, \\forall k \\}$ (e.g., $\\frac{1}{\\sqrt{2R}} [1, -1]^T$ in 1D) to all training signals and find their sparse representations, and then find $\\alpha_{\\mathrm{est}}$ that corresponds to the largest $95 (\\pm 1)\\%$ of non-zero elements of the sparsified training signals. This procedure defines the range $[ \\frac{1}{10} \\alpha_{\\mathrm{est}}, \\alpha_{\\mathrm{est}}]$ to select desirable $\\alpha^\\star$ and its corresponding filter $D^\\star$. We next ran CAOL~\\R{sys:CAOL:orth} with multiple $\\alpha$ values within this range. Selecting $\\{ \\alpha^\\star, D^\\star \\}$ depends on application. \nFor CT MBIR, $D^\\star$ that both has (short) first-order finite difference filters and captures diverse (particularly diagonal) features of training signals, gave good RMSE values and well preserved edges; see Fig.~\\ref{fig:CAOL:filters:fruit&city&ct-ii}(c) and \\cite[Fig.~2]{Chun&Fessler:18Asilomar}. \n\n\n\n\n\n\\begin{figure*}[!t]\n\\centering\n\\small\\addtolength{\\tabcolsep}{-6.5pt}\n\\renewcommand{\\arraystretch}{0.1}\n\n \\begin{tabular}{ccccc}\n\n\n \\small{(a) Ground truth} \n & \n \\small{(b) Filtered back-projection} & \\small{(c) EP} \n & \n \\specialcell[c]{\\small (d) Proposed MBIR \\mbox{\\R{sys:CT&CAOL}}, \\\\ \\small with \\R{eq:autoenc} of $R \\!=\\! K \\!=\\! 25$} \n & \n \\specialcell[c]{\\small (e) Proposed MBIR \\mbox{\\R{sys:CT&CAOL}}, \\\\ \\small with \\R{eq:autoenc} of $R \\!=\\! K \\!=\\! 49$} \\\\\n \n \\begin{tikzpicture}\n \t\t\\begin{scope}[spy using outlines={rectangle,yellow,magnification=1.25,size=15mm,connect spies}]\n \t\t\t\\node {\\includegraphics[viewport=35 105 275 345, clip, width=3.4cm,height=3.4cm]{.\/Fig_CTrecon\/x_true.png}};\n\t\t\t\\spy on (0.8,-0.95) in node [left] at (-0.3,-2.2);\n \t\t\\draw [->, thick, red] (0.1,0.85+0.3) -- (0.1,0.85);\n \t\t\\draw [->, thick, red] (-0.65+0.3,0.1) -- (-0.65,0.1);\n \t\t\\draw [->, thick, red] (-1.1,-1.3) -- (-1.1+0.3,-1.3); \n\t\t\t\\draw [red] (-1.35,0.15) circle [radius=0.28];\n\t\t\t\\draw [red] (-1.35,-0.45) circle [radius=0.28];\n\t\t\\end{scope}\n \\end{tikzpicture} \n &\n \\begin{tikzpicture}\n \t\t\\begin{scope}[spy using outlines={rectangle,yellow,magnification=1.25,size=15mm,connect spies}]\n \t\t\t\\node {\\includegraphics[viewport=35 105 275 345, clip, width=3.4cm,height=3.4cm]{.\/Fig_CTrecon\/x123_fbp.png}};\n\t\t\t\\spy on (0.8,-0.95) in node [left] at (-0.3,-2.2);\n\t\t\t\\node [black] at (0.7,-1.95) {\\small $\\mathrm{RMSE} = 82.8$};\n\t\t\\end{scope}\n \\end{tikzpicture} \n &\n \\begin{tikzpicture}\n \t\t\\begin{scope}[spy using outlines={rectangle,yellow,magnification=1.25,size=15mm,connect spies}]\n \t\t\t\\node {\\includegraphics[viewport=35 105 275 345, clip, width=3.4cm,height=3.4cm]{.\/Fig_CTrecon\/x123_ep.png}};\n\t\t\t\\spy on (0.8,-0.95) in node [left] at (-0.3,-2.2);\n \t\t\\draw [->, thick, red] (0.1,0.85+0.3) -- (0.1,0.85);\n \t\t\\draw [->, thick, red] (-0.65+0.3,0.1) -- (-0.65,0.1);\n \t\t\\draw [->, thick, red] (-1.1,-1.3) -- (-1.1+0.3,-1.3); \n\t\t\t\\node [black] at (0.7,-1.95) {\\small $\\mathrm{RMSE} = 40.8$};\n\t\t\\end{scope}\n \\end{tikzpicture} \n &\n \\begin{tikzpicture}\n \t\t\\begin{scope}[spy using outlines={rectangle,yellow,magnification=1.25,size=15mm,connect spies}]\n \t\t\t\\node {\\includegraphics[viewport=35 105 275 345, clip, width=3.4cm,height=3.4cm]{.\/Fig_CTrecon_v1\/x123_caol5x5_old.png}};\n\t\t\t\\spy on (0.8,-0.95) in node [left] at (-0.3,-2.2);\n \t\t\\draw [->, thick, red] (0.1,0.85+0.3) -- (0.1,0.85);\n \t\t\\draw [->, thick, red] (-0.65+0.3,0.1) -- (-0.65,0.1);\n \t\t\\draw [->, thick, red] (-1.1,-1.3) -- (-1.1+0.3,-1.3); \n\t\t\t\\draw [red] (-1.35,0.15) circle [radius=0.28];\n\t\t\t\\draw [red] (-1.35,-0.45) circle [radius=0.28];\n\t\t\t\\node [blue] at (0.7,-1.95) {\\small $\\mathrm{RMSE} = 35.2$};\n\t\t\\end{scope}\n \\end{tikzpicture} \n &\n \\begin{tikzpicture}\n \t\t\\begin{scope}[spy using outlines={rectangle,yellow,magnification=1.25,size=15mm,connect spies}]\n \t\t\t\\node {\\includegraphics[viewport=35 105 275 345, clip, width=3.4cm,height=3.4cm]{Fig_CTrecon_v1\/x123_caol7x7_old.png}};\n\t\t\t\\spy on (0.8,-0.95) in node [left] at (-0.3,-2.2);\n \t\t\\draw [->, thick, red] (0.1,0.85+0.3) -- (0.1,0.85);\n \t\t\\draw [->, thick, red] (-0.65+0.3,0.1) -- (-0.65,0.1);\n \t\t\\draw [->, thick, red] (-1.1,-1.3) -- (-1.1+0.3,-1.3); \n\t\t\t\\draw [red] (-1.35,0.15) circle [radius=0.28];\n\t\t\t\\draw [red] (-1.35,-0.45) circle [radius=0.28];\n\t\t\t\\node [red] at (0.7,-1.95) {\\small $\\mathrm{RMSE} = 34.7$};\n\t\t\\end{scope}\n \\end{tikzpicture} \n\n \\end{tabular}\n\n \\vspace{-0.25em}\n\\caption{\nComparisons of reconstructed images from different reconstruction methods for sparse-view CT ($123$ views ($12.5$\\% sampling); for the MBIR model \\R{sys:CT&CAOL}, convolutional regularizers were trained by CAOL \\R{sys:CAOL:orth} -- see \\cite[Fig.~2]{Chun&Fessler:18Asilomar}; display window is within $[800, 1200]$ HU) \\cite{Chun&Fessler:18Asilomar}. \nThe MBIR model \\R{sys:CT&CAOL} using convolutional sparsifying regularizers trained via CAOL \\R{sys:CAOL:orth} shows higher image reconstruction accuracy compared to the EP reconstruction; see red arrows and magnified areas. \nFor the MBIR model \\R{sys:CT&CAOL}, the autoencoder (see Remark~\\ref{r:autoenc}) using the filter dimension $R \\!=\\! K \\!=\\! 49$ improves reconstruction accuracy of that using $R \\!=\\! K \\!=\\! 25$; compare the results in (d) and (e).\nIn particular, the larger dimensional filters improve the edge sharpness of reconstructed images; see circled areas. The corresponding error maps are shown in Fig.~\\ref{fig:CTrecon:err} of the supplementary material.\n}\n\\label{fig:CTrecon}\n\\end{figure*}\n\n\n\n\n\n\n\n\\subsection{Sparse-View CT MBIR with Learned Convolutional Sparsifying Regularizer (via CAOL) and BPEG-M} \\label{sec:result:CT}\n\n\n\nIn sparse-view CT using only $12.5$\\% of the full projections views, the CT MBIR \\R{sys:CT&CAOL} using the learned convolutional regularizer via CAOL \\R{sys:CAOL:orth} outperforms EP MBIR; it reduces RMSE by approximately $5.6$--$6.1$HU.\nSee the results in Figs.~\\ref{fig:CTrecon}(c)--(e).\nThe model \\R{sys:CT&CAOL} can better recover high-contrast regions (e.g., bones) -- see red arrows and magnified areas in Fig.~\\ref{fig:CTrecon}(c)--(e). \nNonetheless, the filters with $R\\!=\\!K\\!=\\!5^2$ in the ($\\psi$-weighting) autoencoding CNN, i.e., $\\sum_{k} ( P_f d_k^\\star ) \\circledast \\cH_{\\!\\sqrt{2 \\alpha' \\psi}} ( d_k^\\star \\circledast (\\cdot) )$ in \\R{eq:autoenc}, can blur edges in low-contrast regions (e.g., soft tissues) while removing noise. See Fig.~\\ref{fig:CTrecon}(d) -- the blurry issues were similarly observed in \\cite{Chun&etal:17Fully3D, Zheng&etal:19TCI}. \nThe larger dimensional kernels (i.e., $R\\!=\\!K\\!=\\!7^2$) in the convolutional autoencoder can moderate this issue, while further reducing RMSE values; compare the results in Fig.~\\ref{fig:CTrecon}(d)--(e).\nIn particular, the larger dimensional convolutional kernels capture more diverse features -- see \\cite[Fig.~2]{Chun&Fessler:18Asilomar}) -- and the diverse features captured in kernels are useful to further improve the performance of the proposed MBIR model \\R{sys:CT&CAOL}. (The importance of diverse features in kernels was similarly observed in CT experiments with the learned autoencoders having a fixed kernel dimension; see Fig.~\\ref{fig:CAOL:filters:fruit&city&ct-ii}(c).)\nThe RMSE reduction over EP MBIR is comparable to that of CT MBIR \\R{sys:CT&CAOL} \nusing the $\\{ R, K \\!=\\! 8^2 \\}$-dimensional filters trained via the patch-domain AOL \\cite{Ravishankar&Bressler:15TSP}; \nhowever, at each BPEG-M iteration, this MBIR model using the trained (non-TF) filters via patch-domain AOL \\cite{Ravishankar&Bressler:15TSP} requires more computations than the proposed CT MBIR model \\R{sys:CT&CAOL} using the learned convolutional regularizer via CAOL \\R{sys:CAOL:orth}.\nSee related results and discussion in \nFig.~\\ref{fig:CTrecon:TL} and Section~\\ref{sec:result:supp}, respectively. \n \n\n\n\nOn the algorithmic side, the BPEG-M framework can guarantee the convergence of CT MBIR \\R{sys:CT&CAOL}.\nUnder the sharp majorization regime in BPEG-M, maintaining the majorization sharpness is more critical than having stronger extrapolation effects -- see \\cite[Fig.~3]{Chun&Fessler:18Asilomar}, as similarly shown in CAOL experiments (see Section~\\ref{sec:result:CAOL}).\n\n\n\n\n\n\\section{Conclusion} \\label{sec:conclusion}\n\nDeveloping rapidly converging and memory-efficient CAOL engines is important, since it is a basic element in training CNNs in an unsupervised learning manner (see Appendix~\\ref{sec:CNN}). \nStudying structures of convolutional kernels is another fundamental issue, since it can avoid learning redundant filters or provide energy compaction properties to filters.\nThe proposed BPEG-M-based CAOL framework has several benefits.\nFirst, the orthogonality constraint and diversity promoting regularizer in CAOL are useful in learning filters with diverse structures.\nSecond, the proposed BPEG-M algorithm significantly accelerates CAOL over the state-of-the-art method, BPG \\cite{Xu&Yin:17JSC}, with our sufficiently sharp majorizer designs.\nThird, BPEG-M-based CAOL uses much less memory compared to patch-domain AOL methods \\cite{Yaghoobi&etal:13TSP, Hawe&Kleinsteuber&Diepold:13TIP, Ravishankar&Bressler:15TSP}, and easily allows parallel computing.\nFinally, the learned convolutional regularizer provides the autoencoding CNN architecture in MBIR, and outperforms EP reconstruction in sparse-view CT.\n\n\nSimilar to existing unsupervised synthesis or analysis operator learning methods, the biggest remaining challenge of CAOL is optimizing its model parameters. \nThis would become more challenging when one applies CAOL to train CNNs (see Appendix~\\ref{sec:CNN}).\nOur first future work is developing \\dquotes{task-driven} CAOL that is particularly useful to train thresholding values.\nOther future works include further acceleration of BPEG-M in Algorithm~\\ref{alg:BPGM}, designing sharper majorizers requiring only $O( L R N )$ for the filter update problem of CAOL~\\R{sys:CAOL}, and applying the CNN model learned via \\R{sys:CNN:orth} to MBIR. \n\n\n\n\n\n\n\n\n\n\n\\section*{Appendix}\n\\renewcommand{\\thesubsection}{\\Alph{subsection}}\n\n\\subsection{Training CNN in a unsupervised manner via CAOL} \\label{sec:CNN}\n\nThis section mathematically formulates an unsupervised training cost function for classical CNN (e.g., LeNet-5 \\cite{LeCun&etal:98ProcIEEE} and AlexNet \\cite{Krizhevsky&etal:12NIPS}) and solves the corresponding optimization problem, via the CAOL and BPEG-M frameworks studied in Sections~\\ref{sec:CAOL}--\\ref{sec:CAOL+BPGM}.\nWe model the three core modules of CNN: \\textit{1)} convolution, \\textit{2)} pooling, e.g., average \\cite{LeCun&etal:98ProcIEEE} or max \\cite{Jarrett&etal:09ICCV}, and \\textit{3)} thresholding, e.g., RELU \\cite{Nair&Hinton:10ICML}, while considering the TF filter condition in Proposition~\\ref{p:TFconst}. \nParticularly, the orthogonality constraint in CAOL \\R{sys:CAOL:orth} leads to a sharp majorizer, and BPEG-M is useful to train CNNs with convergence guarantees.\nNote that it is unclear how to train such diverse (or incoherent) filters described in Section~\\ref{sec:CAOL} by the most common CNN optimization method, the stochastic gradient method in which gradients are computed by back-propagation. The major challenges include \\textit{a)} the non-differentiable hard thresholding operator related to $\\ell^0$-norm in \\R{sys:CAOL}, \\textit{b)} the nonconvex filter constraints in \\R{sys:CAOL:orth} and \\R{sys:CAOL:div}, \\textit{c)} using the identical filters in both encoder and decoder (e.g., $W$ and $W^H$ in Section~\\ref{sec:prf:p:TF}), and \\textit{d)} vanishing gradients.\n\nFor simplicity, we consider a two-layer CNN with a single training image, but one can extend the CNN model \\R{sys:CNN:orth} (see below) to \\dquotes{deep} layers with multiple images. The first layer consists of \\textit{1c)} convolutional, \\textit{1t)} thresholding, and \\textit{1p)} pooling layers; the second layer consists of \\textit{2c)} convolutional and \\textit{2t)} thresholding layers.\nExtending CAOL \\R{sys:CAOL:orth}, we model two-layer CNN training as the following optimization problem:\n\\begingroup\n\\setlength{\\thinmuskip}{1.5mu}\n\\setlength{\\medmuskip}{2mu plus 1mu minus 2mu}\n\\setlength{\\thickmuskip}{2.5mu plus 2.5mu}\n\\allowdisplaybreaks\n\\ea{\n\\label{sys:CNN:orth}\n\\argmin_{\\{ d_{k}^{[1]}, d_{k,k'}^{[2]} \\}} \\min_{\\{ z_{k}^{[1]}, z_{k'}^{[2]} \\}}\n&~ \\sum_{k=1}^{K_1} \\frac{1}{2} \\left\\| d_{k}^{[1]} \\circledast x - z_{k}^{[1]} \\right\\|_2^2 + \\alpha_1 \\left\\| z_{k}^{[1]} \\right\\|_0 \\nn\n\\\\\n+ &\\, \\frac{1}{2} \\left\\| \\left( \\sum_{k=1}^{K_1} \\left[ \\arraycolsep=2pt \\begin{array}{c} d_{k,1}^{[2]} \\circledast P z_{k}^{[1]} \\\\ \\vdots \\\\ d_{k,K_2}^{^{[2]}} \\circledast P z_{k}^{[1]} \\end{array} \\right] \\right) - \\left[ \\arraycolsep=2pt \\begin{array}{c} z_{1}^{[2]} \\\\ \\vdots \\\\ z_{K_2}^{[2]} \\end{array} \\right] \\right\\|_2^2 \\nn\n\\\\\n+ &\\, \\alpha_2 \\sum_{k'=1}^{K_2} \\left\\| z_{k'}^{[2]} \\right\\|_0 \\nn\n\\\\\n\\mathrm{subject~to} ~~~~&~ D^{[1]} \\big( D^{[1]} \\big)^H = \\frac{1}{R_1} \\cdot I, \\nn\n\\\\\n&~ D_k^{[2]} \\big( D_k^{[2]} \\big)^H = \\frac{1}{R_2} \\cdot I, \\quad k = 1,\\ldots,K_1, \n\\tag{A1}\n} \n\\endgroup\nwhere $x \\in \\bbR^{N}$ is the training data, $\\{ d_{k}^{[1]} \\in \\bbR^{R_1} : k = 1,\\ldots,K_1 \\}$ is a set of filters in the first convolutional layer, $\\{ z_{k}^{[1]} \\in \\bbR^{N} : k = 1,\\ldots,K_1\\}$ is a set of features after the first thresholding layer, $\\{ d_{k,k'}^{[2]} \\in \\bbR^{R_2} : k' = 1,\\ldots,K_2 \\}$ is a set of filters for each of $\\{ z_{k}^{[1]} \\}$ in the second convolutional layer, $\\{ z_{k'}^{[2]} \\in \\bbR^{N\/\\omega} : k = 1,\\ldots,K_2\\}$ is a set of features after the second thresholding layer, $D^{[1]}$ and $\\{ D_k^{[2]} \\}$ are similarly given as in \\R{eq:D}, $P \\in \\bbR^{N\/\\omega \\times \\omega}$ denotes an average pooling \\cite{LeCun&etal:98ProcIEEE} operator (see its definition below), and $\\omega$ is the size of pooling window. \nThe superscripted number in the bracket of vectors and matrices denotes the $(\\cdot)\\rth$ layer.\nHere, we model a simple average pooling operator $P \\in \\bbR^{(N\/\\omega) \\times \\omega}$ by a block diagonal matrix with row vector $\\frac{1}{\\omega} 1_\\omega^T \\in \\bbR^{\\omega}$: $P := \\frac{1}{\\omega} \\bigoplus_{j=1}^{N\/\\omega} 1_\\omega^T$.\nWe obtain a majorization matrix of $P^T P$ by $P^T P \\preceq \\diag( P^T P 1_N ) = \\frac{1}{\\omega} I_N$ (using Lemma~\\ref{l:diag(|At|W|A|1)}).\nFor 2D case, the structure of $P$ changes, but $P^T P \\preceq \\frac{1}{\\omega} I_N$ holds.\n\n\n\n\n\n\n\nWe solve the CNN training model in \\R{sys:CNN:orth} via the BPEG-M techniques in Section~\\ref{sec:CAOL+BPGM}, and relate the solutions of \\R{sys:CNN:orth} and modules in the two-layer CNN training. The symbols in the following items denote the CNN modules.\n\\begin{itemize}\n[\\setlength{\\IEEElabelindent}{\\IEEEilabelindentA}]\n\\item[\\textit{1c})] Filters in the first layer, $\\{ d_{k}^{[1]} \\}$: Updating the filters is straightforward via the techniques in Section~\\ref{sec:prox:filter:orth}. \n\n\\item[\\textit{1t})] Features at the first layers, $\\{ z_{k}^{[1]} \\}$: Using BPEG-M with the $k\\rth$ set of TF filters $\\{ d_{k,k'}^{[2]} : k' \\}$ and $P^T P \\preceq \\frac{1}{\\omega} I_N$ (see above), the proximal mapping for $z_{k}^{[1]}$ is\n\\begingroup\n\\setlength{\\thinmuskip}{1.5mu}\n\\setlength{\\medmuskip}{2mu plus 1mu minus 2mu}\n\\setlength{\\thickmuskip}{2.5mu plus 2.5mu}\n\\be{\n\\label{eq:CNN:layer1:spCd}\n\\min_{z_{k}^{[1]}} \\frac{1}{2} \\left\\| d_{k}^{[1]} \\circledast x - z_{k}^{[1]} \\right\\|_2^2 + \\frac{1}{2\\omega'} \\left\\| z_{k}^{[1]} - \\zeta_{k}^{[k]} \\right\\|_2^2 + \\alpha_1 \\left\\| z_{k}^{[1]} \\right\\|_0,\n}\n\\endgroup\nwhere $\\omega' = \\omega \/ \\lambda_Z$ and $\\zeta_{k}^{[k]}$ is given by \\R{update:x}. Combining the first two quadratic terms in \\R{eq:CNN:layer1:spCd} into a single quadratic term leads to an optimal update for \\R{eq:CNN:layer1:spCd}:\n\\bes{\nz_{k}^{[1]} = \\cH_{\\! \\sqrt{ 2 \\frac{\\omega' \\alpha_1}{\\omega' + 1} }} \\left( d_{k}^{[1]} \\circledast x + \\frac{1}{\\omega'} \\zeta_{k}^{[k]} \\right), \\quad k \\in [K],\n} \nwhere the hard thresholding operator $\\cH_a (\\cdot)$ with a thresholding parameter $a$ is defined in \\R{eq:def:hardthr}.\n\n\n\\item[\\textit{1p})] Pooling, $P$: Applying the pooling operator $P$ to $\\{ z_{k}^{[1]} \\}$ gives input data -- $\\{ P z_{k}^{[1]} \\}$ -- to the second layer.\n\n\\item[\\textit{2c})] Filters in the second layer, $\\{ d_{k,k'}^{[2]} \\}$: We update the $k\\rth$ set filters $\\{ d_{k,k'}^{[2]} : \\forall k' \\}$ in a sequential way. Updating the $k\\rth$ set filters is straightforward via the techniques in Section~\\ref{sec:prox:filter:orth}. \n\n\\item[\\textit{2t})] Features at the second layers, $\\{ z_{k'}^{[2]} \\}$: The corresponding update is given by\n\\bes{\nz_{k'}^{[2]} = \\cH_{\\!\\sqrt{2 \\alpha_2}} \\left( \\sum_{k=1}^{K_1} d_{k,k'}^{[1]} \\circledast P z_{k}^{[1]} \\right), \\quad k' \\in [K_2].\n}\n\\end{itemize}\n\nConsidering the introduced mathematical formulation of training CNNs \\cite{LeCun&etal:98ProcIEEE} via CAOL, BPEG-M-based CAOL has potential to be a basic engine to rapidly train CNNs with big data (i.e., training data consisting of many (high-dimensional) signals).\n\n\n\n\\subsection{Examples of $\\{ f(x;y), \\cX \\}$ in MBIR model \\R{eq:mbir:learn} using learned regularizers} \\label{sec:egs}\n\n\nThis section introduces some potential applications \nof using MBIR model \\R{eq:mbir:learn} using learned regularizers\nin imaging processing, imaging, and computer vision.\nWe first consider quadratic data fidelity function in the form of $f(x;y) = \\frac{1}{2} \\| y - A x \\|_W^2$.\nExamples include\n\\begin{itemize}\n\\item Image debluring (with $W \\!=\\! I$ for simplicity), \nwhere $y$ is a blurred image, $A$ is a blurring operator, and $\\cX$ is a box constraint; \n\\item Image denoising (with $A \\!=\\! I$),\nwhere $y$ is a noisy image corrupted by additive white Gaussian noise (AWGN),\n$W$ is the inverse covariance matrix corresponding to AWGN statistics,\nand $\\cX$ is a box constraint;\n\\item Compressed sensing (with $\\{ W \\!=\\! I, \\cX \\!\\in\\! \\bbC^{N'} \\}$ for simplicity) \\cite{Chun&Adcock:17TIT, Chun&Adcock:18ACHA},\nwhere $y$ is a measurement vector, \nand $A$ is a compressed sensing operator,\ne.g., subgaussian random matrix, bounded orthonormal system, subsampled isometries, certain types of random convolutions;\n\\item Image inpainting (with $W \\!=\\! I$ for simplicity),\nwhere $y$ is an image with missing entries, $A$ is a masking operator, and $\\cX$ is a box constraint;\n\\item Light-field photography from focal stack data \nwith $f(x;y) = \\sum_{c} \\| y_c - \\sum_{s} A_{c,s} x_{s} \\|_2^2$,\nwhere $y_c$ denotes measurements collected at the $c\\rth$ sensor, \n$A_{c,s}$ models camera imaging geometry at the $s\\rth$ angular position for the $c\\rth$ detector, \n$x_s$ denotes the $s\\rth$ sub-aperture image, $\\forall c,s$,\nand $\\cX$ is a box constraint \\cite{Block&Chun&Fessler:18IVMSP, Chun&etal:18arXiv:momnet}.\n\\end{itemize}\nExamples that use nonlinear data fidelity function include \nimage classification using the logistic function \\cite{Mairal&etal:09NIPS}, \nmagnetic resonance imaging considering unknown magnetic field variation \\cite{Fessler:10SPM}, \nand positron emission tomography \\cite{Lim&etal:19TMI}.\n\n\n\n\n\\subsection{Notation} \\label{sec:notation}\n\nWe use $\\nm{\\cdot}_{p}$ to denote the $\\ell^p$-norm and write $\\ip{\\cdot}{\\cdot}$ for the standard inner product on $\\bbC^N$. \nThe weighted $\\ell^2$-norm with a Hermitian positive definite matrix $A$ is denoted by $\\nm{\\cdot}_{A} = \\nm{ A^{1\/2} (\\cdot) }_2$.\n$\\nm{\\cdot}_{0}$ denotes the $\\ell^0$-quasi-norm, i.e., the number of nonzeros of a vector. \nThe Frobenius norm of a matrix is denoted by $\\| \\cdot \\|_{\\mathrm{F}}$. \n$( \\cdot )^T$, $( \\cdot )^H$, and $( {\\cdot} )^*$ indicate the transpose, complex conjugate transpose (Hermitian transpose), and complex conjugate, respectively. \n$\\diag(\\cdot)$ denotes the conversion of a vector into a diagonal matrix or diagonal elements of a matrix into a vector.\n$\\bigoplus$ denotes the matrix direct sum of matrices.\n$[C]$ denotes the set $\\{1,2,\\ldots,C\\}$. \nDistinct from the index $i$, we denote the imaginary unit $\\sqrt{-1}$ by $\\imath$. \nFor (self-adjoint) matrices $A,B \\in \\bbC^{N \\times N}$, the notation $B \\preceq A$ denotes that $A-B$ is a positive semi-definite matrix. \n\n\n\n\n\\section*{Acknowledgment}\n\nWe thank Xuehang Zheng for providing CT imaging simulation setup, and Dr. Jonghoon Jin for constructive feedback on CNNs.\n\n\n\n\n\\bibliographystyle{IEEEtran}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\\label{sec:introduction}\n\nThe Planck 2013 and 2015 data releases open new directions in precision cosmology with regard to a more advanced\ninvestigation of the statistical isotropy and non-Gaussianity of the cosmic microwave background (CMB)~\\cite{Planck1, Planck2, Planck_is, Planck_is2015}.\nWhile generally confirming the Gaussianity and the statistical isotropy of the CMB (for the relevant multipole domain $\\ell\\ge 50$), the Planck\nscience team confirmed the existence of a variety of anomalies in the temperature anisotropy on large angular scales ($\\ell\\le 50$)\npreviously seen in WMAP data~\\cite{Planck_is,Planck_is2015}. Among these are the lack of power in the quadrupole component~\\cite{low_quadrupole} (see, however,~\\cite{wmap9b}\nand~\\cite{Planck_is}), the alignment of the quadrupole and octupole components~\\cite{Copi1, wmap9b, Planck_is}, the unusual symmetry of the\noctupole~\\cite{mhansen}, anisotropies in the temperature angular power spectrum~\\cite{WMAP7:powerspectra, Planck_is, Hansen2004, Eriksen2003},\n preferred directions~\\cite{Pref.Direction1, Pref.Direction2, Copi2, Planck_is, Planck_is2015, Akrami2014}, asymmetry in the power of even and odd modes~\\cite{parity1, parity2, wmap9b, Planck_is, Planck_is2015} and the Cold Spot~\\cite{vielva2003, Coldspot1, Planck_is, Planck_is2015}.\n Some of these anomalies are probably a consequence of the residuals of foreground effects that could be a major source of contamination in the primordial $E$- and $B$-modes of polarization (in this connection see~\\cite{Planck_bicep}).\n\nThe statistics of $B$-mode polarization that can be derived from ongoing and planned CMB experiments will be crucial for the determination of the\ncosmological gravitational waves associated with inflation~\\cite{inflation} at the range of multipoles $50\\le \\ell\\le 150$, closer to\nthe domain of interest for BICEP2 and Planck~\\cite{bicep,Planck_bicep}. It seems likely that $B$-mode polarization in this range is affected\nby Galactic dust emission, the statistical properties of which are very poorly known. (For $\\ell >150$ we expect contamination of the $B$-modes\ndue to lensing effects the precise nature of which is also not fully understood.) In the absence of such knowledge it is difficult to make \\emph{a priori}\nproposals for the best estimator of non-Gaussianity and statistical anisotropies in the derived $B$-modes due to possible contamination by foreground residuals.\nThus, we believe that it is of value to propose additional model-independent tests aimed at providing an improved quantitative understanding of the magnitude of non-Gaussianity in current CMB data. Such tests would also be useful for the analysis of forthcoming CMB data sets. In this paper\nwe propose use of the Kullback--Leibler (KL) divergence as such a test. The goal of this paper is to illustrate the utility of the KL divergence in studying the properties of the CMB signal --- a Gaussian or almost Gaussian signal. The KL divergence is likely to be even more useful for very non-Gaussian cases such as the statistical behaviour of the Minkowski functionals for a single map or the pixel-pixel cross-correlation coefficient between two maps, when calculated in small areas. We will consider such issues in a separate publication.\n\nThe one-point probability distribution function (PDF) would seem to be a reasonable starting point for the investigation of non-Gaussianity.\nSuch tests have been applied by the Planck team to a variety of derived CMB temperature maps including the SMICA, NILC, SEVEM and\nCommander maps~\\cite{Planck_is,Planck_is2015}. In practice, these tests involve comparison of the various temperature fluctuation maps with an ensemble\nof simulated maps. The CMB temperature is characterized by a power spectrum, $C_\\ell^\\text{Planck}$, which\ncorresponds to the Planck 2013 concordance \\LCDM\\ model with the cosmological parameters listed in~\\cite{Planck_cp}.\nThe simulated maps are obtained as Monte Carlo (MC) Gaussian draws on this power spectrum (in harmonic space). When processed with the Planck component separation tool, the resulting simulation maps contain both CMB information as well as various residuals from foregrounds, the uncertainties of instrumental effects, etc. For the Planck 2013 data release, the corresponding $10^3$ full focal plane simulations are referred to as the FFP7 simulations.\nThey reflect the intrinsic properties of the SMICA, NILC, SEVEM and Commander maps. Differences between the FFP7 maps and the various empirical maps can provide useful information regarding non-Gaussianity. In the following, we shall restrict our attention to the SMICA map.\n\nIn practice, the non-parametric Kolmogorov--Smirnov (KS) test is often used to assess the similarity of two distributions. The KS test\ncharacterizes the difference between the two cumulative distribution functions (CDF) in terms of the maximum absolute deviation between them. The KS estimator, $\\kappa$, is defined as\n\\begin{equation}\n\\kappa= \\sqrt{n} \\max[|F(x)-F_n(x)|] \\, ,\n\\label{eq:KS_definition}\n\\end{equation}\nwhere $F(x)$ is the theoretical expectation of the CDF and $F_n(x)$ is obtained from a data sample with $n$ elements. Here, both $F(x)$ and $F_n(x)$ must be normalized to the range $[0,1]$. Note that $F(x)$ is normally a continuous function or should at least be defined for all possible values of the data sample $F_n(x)$. It is clear from Eq.~\\eqref{eq:KS_definition} that the KS estimator $\\kappa$ is local in the sense that its value will be determined at a point, $x$, where the PDFs corresponding to $F(x)$ and $F_n(x)$ cross. We note that the use of PDFs in Eq.~\\eqref{eq:KS_definition} instead of CDFs would result in a\nmaximal sensitivity to the largest local anomaly. \n\nUnlike the case of vectors (where the scalar product provides a standard measure), there is no generic ``best''\nmeasure for quantifying the similarity of two distributions. Thus, we believe that it is also useful to consider the Kullback--Leibler divergence for two discrete probability distributions, $P$ and $Q$. The KL divergence on the set of points $i$ is defined~\\cite{kl} as\n\\be \\label{eq:KL_definition}\nK(P\\|Q)=\\sum_i P_i\\log\\left(\\frac{P_i}{Q_i}\\right) \\, .\n\\ee\nIn other words, the KL divergence is the expectation value of the logarithmic difference between the two probability distributions as computed with\nweights of $P_i$. Typically, $P$ represents the distribution of the data, while $Q$ represents a theoretical expectation of the data. Unlike the\nKS test, the KL divergence is non-local. Indeed, we shall indicate below that it is in a sense ``maximally'' non-local. It is familiar in information\ntheory, where it represents the difference between the intrinsic entropy,\n$H_P$, of the distribution $P$ and the cross-entropy, $H_{PQ}$, between $P$ and $Q$,\n\\be\nK(P\\|Q)=H_{PQ}-H_P, \\qquad H_P=-\\sum_i P_i\\log P_i, \\qquad H_{PQ}=-\\sum_i P_i\\log Q_i \\, .\n\\label{eq:entropy}\n\\ee\nIn more practical terms, consider the most probable result of $N$ independent random draws on the distribution $P$. When $N$ is large, the number of\ndraws at point $i$ is simply $n_i = N P_i$. Now construct the probabilities, $\\Pi_P$ and $\\Pi_Q$, that this most probable result was drawn at random\non distribution $P$ or $Q$, respectively. The KL divergence of Eq.~\\eqref{eq:KL_definition} is simply $N^{-1}\\log{(\\Pi_P \/ \\Pi_Q)}$.\n We note that simulations of the CMB\nmap, drawn independently in harmonic space, have correlations in pixel space.\n Nevertheless we regard this argument as motivation for applying the KL divergence to CMB pixels.\n\nThe main goal of this paper is to illustrate the implementation of the KL divergence for\nthe analysis of the statistical properties of the derived CMB maps in the low multipole domain as a complementary test to the methods listed in~\\cite{Planck_is,Planck_is2015}. The structure of the paper is as follows. In Section~\\ref{sec:KL_divergence_and_Planck_data} we present some properties of the KL divergence. We also use it to analyze the Planck SMICA map and compare it to both the FFP7 set and to a purely Gaussian ensemble. In addition, we compare the two ensembles and compare the KL divergence to the KS test in the low multipole domain of the CMB map. In Section~\\ref{sec:discussion} we discuss the results. Note that the Planck papers~\\cite{Planck_is,Planck_is2015} test the Gaussianity of the one-point PDF by analyzing its variance, skewness and kurtosis. In this sense, the KL divergence is simply another test on the global shape of the PDF. Here, we restrict our analysis by using the SMICA map and the corresponding simulations. The extension of the method to $E$- and $B$-modes of polarization does not require any modification.\n\n\n\n\\section{KL divergence and Planck data}\n\\label{sec:KL_divergence_and_Planck_data}\n\n\\subsection{Preliminary remarks: Properties of the KL divergence}\n\\label{subsec:Preliminary_remarks_KL_divergence}\n\nAs noted above, the KL divergence provides a measure of the similarity of two known distributions. In many cases, however, one of these\ndistributions is not known and must rather be approximated by the average PDF for a statistical ensemble of realizations of the random field. The\nquestion then arises of how closely the resulting proxy reflects the properties of the true underlying distribution. To offer some insight in this\nmatter, we consider a toy model based on a discrete Gaussian distribution, $P_k$, with\n\\be\nP_k \\sim \\exp{\\left[ - k^2\/25 \\right]}\\, ,\n\\ee\nand $k$ an integer between $-10$ and $+10$ subject to the obvious normalization condition. The mean value of\n$k$ is $0$ and the variance is approximately $25\/2$. Suppose for simplicity that we define an individual data\nset as $N$ random draws on this distribution. Each such data set can be regarded as a proxy, $Q_i$ for the underlying\ndistribution, and can be used to calculate the KL divergence $K(P\\|Q)$ defined in Eq.~\\eqref{eq:KL_definition}. For a given value of $N$, we repeat this\nprocess $M$ times and compute the average KL divergence, $\\overline{K}$, and the root-mean-square (RMS) deviation of the KL divergence, $\\Delta K = \\sqrt{\\overline{K^2} - \\overline{K}{}^2 }$. The results of\nthis exercise are shown in Table~\\ref{tab:Change_N}, where we have used the common value of $M=1000$.\n\\begin{table}\n\t\\centering\n\t\\caption{Mean and RMS values of the KL divergence for a discrete Gaussian distribution.}\n\t\\vspace{10pt}\n\t\\begin{tabular}{rcccc}\n\t\t\\hline\n\t\t\\noalign{\\smallskip}\n\t\t\\multicolumn{1}{c}{$N$} & $\\overline{K}$ & $\\Delta K$ & $N\\overline{K}$ & $N\\Delta K$ \\\\\n\t\t\\hline\n\t\t 4000 & 0.002540 & 0.000796 & 10.160 & 3.18 \\\\\n\t\t 8000 & 0.001253 & 0.000413 & 10.024 & 3.30 \\\\\n\t\t16000 & 0.000636 & 0.000200 & 10.176 & 3.20 \\\\\n\t\t\\hline\n\t\\end{tabular}\n\t\\label{tab:Change_N}\n\\end{table}\nSeveral things seem clear. Both the average value of the KL divergence and the RMS deviation from this average value vanish like $1\/N$ for large $N$. From general arguments,\nthe KL divergence cannot be negative. Fig.~\\ref{fig:KL_hist_change_N} shows a histogram of the distribution of KL divergences\nobtained with $M = 20000$ for the cases $N=100$ and $N=1000$.\n\\begin{figure}\n\t\\centering\n\t\\includegraphics[width=\\halfW]{KL_hist_change_N.ps}\n\t\\caption{Histogram of KL divergences with $M=20000$ for $N=100$ (\\emph{black})\n\tand $N=1000$ (\\emph{red}). Note that the horizontal axis is measured in units of the\n\tcorresponding value of $\\overline{K}$.}\n\t\\label{fig:KL_hist_change_N}\n\\end{figure}\nThe KL divergences here are measured\nin units of the corresponding value of $\\overline{K}$. The fact that these distributions scale like $1\/N$ is\nobvious. These histograms suggest power law suppression near zero and Gaussian behaviour\nfor large values of the KL divergence.\n\n\nThe results shown in Fig.~\\ref{fig:KL_hist_change_N} are not specific to the KL divergence, and qualitatively\nsimilar results would be obtained for any measure chosen to describe the similarity of two distributions.\nAll that is required is that the measure chosen is always positive and vanishes when the distributions\nbeing compared are identical. When drawn as here, each individual data set can be thought of\nas a combination of the ``exact'' distribution plus the amplitudes of $(N-1)$ ``fluctuations''.\\footnote{Due to\nthe fact that each data set contains exactly $N$ draws.} Obviously, the amplitudes of every one of these\nfluctuations must be exactly zero if the measure is to be $K = 0$. Of course, there are many combinations of the\nfluctuation amplitudes that will give any fixed non-zero value of the measure, and their number increases as\n$K$ grows from zero. In contrast, for the case of only two options (e.g., ``heads'' and ``tails'') subject to a constraint\non the total number of draws, there is only a single degree of freedom. In this case $K = 0$ is actually the most\nprobable value.\n\nA few additional remarks can help clarify the properties of the KL divergence. Consider that each individual\ndata set is sufficiently large that $Q_i = P_i + \\delta P_i$ where $\\delta P_i$ is small. Under these conditions,\nterms linear in $\\delta P_i$ vanish as a consequence of normalization, and the KL divergence is given simply as\n\\be\nK = \\frac{1}{2} \\, \\sum_i \\, \\frac{\\delta P_i^2}{P_i} \\, .\n\\label{eq:KL_small_fluctuation}\n\\ee\nWe see that the distributions $P$ and $Q$ are now treated symmetrically in spite of the asymmetry that is\napparent in Eq.~\\eqref{eq:KL_definition}. Elementary arguments suggest that the average number of draws on bin $i$ will be\n$N P_i \\pm \\sqrt{N P_i}$. The corresponding proxy for the underlying distribution will be $P_i \\pm \\sqrt{P_i \/ N}$.\nGiven this result, Eq.~\\eqref{eq:KL_small_fluctuation} suggests that, for fixed bin sizes and in the limit $N \\to \\infty$, each bin will make\na contribution to the KL divergence of roughly equal size. It is hard to imagine a greater degree of non-locality. Moreover, in this limit the KL divergence is expected to be of order $N_b\/N$, where $N_b$ is the number of bins. In other words, if we define\n\\be \\label{eq:alpha_definition}\n\\alpha=\\frac{N}{N_b}K\\, ,\n\\ee\nwe expect that $\\alpha\\sim\\mathcal{O}(1)$. This realization allows us to assign a rough scale for the expected KL divergence. If two given distributions yield a much larger value of $\\alpha$, we can conclude that they differ significantly from each other without resorting to comparisons with an ensemble. In the example shown in Table~\\ref{tab:Change_N} we are using $N_b=30$, meaning $\\alpha$ is small. This is to be expected since we are indeed using the true distribution to draw the data under examination.\n\n\n\\subsection{Preliminary remarks: Planck data}\n\\label{subsec:Preliminary_remarks_Planck_data}\n\nWe have performed the KL divergence test on the CMB map obtained by the Planck collaboration~\\cite{Planck_cp} using the SMICA method. Since we are interested in the\nstatistical properties of the CMB map on large scales, we first degrade the map from its native \\textsc{HEALPix}~\\cite{HEALPix} resolution of $\\Nside=2048$\nto $\\Nside=32$. We then construct the convolution of this map with a Gaussian smoothing kernel of $5\\deg$ FWHM and retain only the harmonic coefficients with\n$\\ell \\le \\lmax=96$. The SMICA map provides a useful estimation of the CMB temperature fluctuations for a very large fraction --- but not all --- of the sky. We use the SMICA inpainting mask to exclude heavily contaminated regions, mainly the Galactic plane. At the resolution considered here, the mask removes about 6\\% of the pixels, leaving the number of pixels under consideration to be $N=11565$. Since the analysis is performed in pixel space,\napplication of the mask is trivial.\n\nIn order to estimate the statistical significance of our results, we compare them to ensembles of realizations. In this work we use two different ensembles.\nThis is done in order to cross-check the significance estimations and also allows us to compare the two ensembles. The first of these ensembles is the\nFFP7 set described above. We degrade and smooth the FFP7 maps in the same manner as we did the SMICA map. We expect that the effects of detector noise will be minor on the large scales considered here. To test this expectation we therefore also make use of the best-fit power spectrum,\n$C_\\ell^\\text{Planck}$~\\cite{Planck_cp}, to generate an ensemble of $10^3$ Gaussian random realizations free of residuals. As in the\ncase of the FFP7 ensemble, we restrict the multipole domain to $\\lmax=96$ and smooth the harmonic coefficients with a Gaussian filter of $5\\deg$ FWHM.\nWe also multiply the coefficients with the pixel window function associated with an $\\Nside=32$ pixelization before converting them to\nan $\\Nside=32$ map.\n\nIn our analysis we calculate the KL divergence for the SMICA map in pixel space. We calculate a histogram of the temperature fluctuations in the\nunmasked pixels by taking bins of width $8\\uK$ in the range $[-200,200]\\uK$, meaning that the number of bins is $N_b=51$. Values outside this range are attributed to the edge bins. This histogram is taken as the $P$ probability distribution of Eq.~\\eqref{eq:KL_definition}. For $Q$, the expected distribution, we turn to the ensemble of simulations, either FFP7 or\nthe Gaussian realizations. We calculate the histogram for each simulation (using the same range and binning), and take $Q$ to be the mean of all histograms. These histograms are shown in Fig.~\\ref{subfig:maps_histograms} together with error bars showing the 5--95\\% range for the FFP7 set.\n\\begin{figure}\n\t\\centering\n\t\\subfigure[]{\n\t\t\\includegraphics[width=\\halfW]{T_hist.ps}\n\t\t\\label{subfig:maps_histograms}\n\t}\n\t\\subfigure[]{\n\t\t\\includegraphics[width=\\halfW]{KL_hist.ps}\n\t\t\\label{subfig:KL_histograms}}\n\t \\caption{\\subref{subfig:maps_histograms}~Number of counts versus amplitude of the SMICA map (\\emph{red}), the FFP7 (\\emph{black}) and\n\tthe Gaussian (\\emph{blue}) ensembles. The error bars show the 5--95\\% range for the FFP7 set. \\subref{subfig:KL_histograms}~Histograms of normalized KL\n\tdivergence values for the FFP7 ensemble (\\emph{black}) and the Gaussian ensemble (\\emph{blue}). The values for the SMICA map are shown as red vertical\n\tlines, compared to the FFP7 mean distribution (\\emph{solid}) and to the Gaussian mean distribution (\\emph{dotted}).}\n\t\\label{fig:KLResults}\n\\end{figure}\nIt appears that the histogram of the SMICA map deviates from the reference histograms by~$\\approx2\\sigma$ primarily in the vicinity of the peak of the distribution. However, this estimation relies on a local feature. The KL divergence provides us with a recipe to sum all the deviations from the entire range of the distribution with appropriate weights.\n\nBefore using Eq.~\\eqref{eq:KL_definition} to calculate the KL divergence, it is necessary to pay particular attention to bins in which either $P$ or $Q$ has small values.\nThe case in which $P_i=0$ is not problematic since $P_i\\log P_i \\to 0$ in this limit. Bins for which $Q_i=0$, however, should not be included since the\nKL divergence is logarithmically divergent as $Q_i \\to 0$. Such a result is not unreasonable since it is impossible to draw to a bin if its probability is\nstrictly $0$. In practice, however, small values of $Q_i$ are merely a consequence of the size of our ensemble. We have chosen to ignore bins for\nwhich $Q_i<5$~pixels in order to minimize the sensitivity to small non-statistical fluctuations in the extreme tails of the $Q$ distribution.\n\n\n\\subsection{The Basic Results}\n\\label{subsec:The_Basic_Results}\n\nWe have calculated the KL divergences between the SMICA histogram and the histograms made from the FFP7 and the Gaussian ensembles. After normalization using the number of valid pixels, $N$, and the number of bins, $N_b$, in Eq.~\\eqref{eq:alpha_definition}, we have obtained $\\alpha=6.47$ and $6.28$, respectively. If these values were significantly larger than the expected order of magnitude, we would conclude that the distributions were in disagreement. Since this is not the case, we must compare the results to ensembles of values of $\\alpha$.\nIn order to calculate the $p$-values, we repeat the calculation of the KL divergence, replacing the distribution of the map, $P$, with that\nof each of the random simulations. This results in two histograms of normalized $K$ values, i.e.\\ $\\alpha$ values, for FFP7 (Gaussian) maps compared to the FFP7 (Gaussian) mean distribution,\nshown in Fig.~\\ref{subfig:KL_histograms}. It is evident that, as expected, $\\alpha\\lesssim10$ for most simulations. We find that $5.6\\%$ of the FFP7 simulations and $6.3\\%$ of the Gaussian simulations get a higher KL\ndivergence than the SMICA map. We see that the KL divergence of the SMICA map from the expected distribution is not significant. As expected, differences\nbetween the two reference ensembles, the FFP7 simulations and the pure Gaussian realizations, are quite small.\nIn order to demonstrate explicitly the similarity between the two ensembles with respect to the KL divergence, we have also tested each of the FFP7 simulations\nagainst the mean distribution of the Gaussian ensemble and vice versa. The results of this calculation are extremely similar to those shown in Fig.~\\ref{subfig:KL_histograms}, and we can conclude that the added complexity of the FFP7 simulations relative to that of simple Gaussian realizations plays a minor role at this resolution.\n\n\nIn addition to the KL divergence, and as a basis for comparison, we also use the KS test to compare between the histogram of the SMICA map and the mean histogram for each of the ensembles. The KS test, defined in Eq.~\\eqref{eq:KS_definition}, requires the use of the CDF. For the SMICA map, the CDF is calculated from the data without any binning. The reference CDF, however, is calculated by first fitting the mean histogram of the ensemble (either FFP7 or the Gaussian realizations) to a Gaussian, and then using the fitted parameters in the expression for a Gaussian CDF. As in the case of the KL divergence, for each ensemble, we compare the SMICA map to the mean histogram of the ensemble and also create a histogram of KS test values by taking each realization separately and comparing it to the mean. The results are shown in Fig.~\\ref{fig:KSResults}.\n\\begin{figure}\n\t\\centering\n\t\\includegraphics[width=\\halfW]{KS_hist.ps}\n\t\\caption{Histograms of KS test values for the FFP7 ensemble (\\emph{black}) and the Gaussian ensemble (\\emph{blue}). The values for the SMICA map are shown as red vertical\nlines, compared to the FFP7 mean distribution (\\emph{solid}) and to the Gaussian mean distribution (\\emph{dotted}).}\n\t\\label{fig:KSResults}\n\\end{figure}\nThe KS test values we get when comparing the SMICA map to the FFP7 and Gaussian simulations are $\\kappa=8.32$ and $8.21$, respectively. The corresponding $p$-values are $3.0\\%$ and $2.6\\%$. Again we see that the results for the two ensembles are in good agreement. Moreover, while the $p$-values of the KS test are lower than those of the KL divergence, the SMICA map still appears to be consistent with the reference ensembles and not anomalous.\n\n\nAs we can see from Fig.~\\ref{subfig:maps_histograms}, there are well-defined temperature ranges in which the SMICA histogram is above or\nbelow the reference. Thus, in Fig.~\\ref{fig:smica_regions} we plot the SMICA map, showing only the temperature range $|T|\\le50\\uK$ where the SMICA\nhistogram is above the reference and the temperature range $50\\uK\\le|T|\\le120\\uK$ where it is below.\n\\begin{figure}\n\t\\centering\n\t\\subfigure[]{\n\t\t\\includegraphics[width=\\thirdW,angle=90]{Map50.ps}\n\t\t\\label{subfig:smica_-50_50}\n\t}\n\t\\subfigure[]{\n\t\t\\includegraphics[width=\\thirdW,angle=90]{Map120.ps}\n\t\t\\label{subfig:smica_abs_50_120}\n\t}\n\t\\caption{The SMICA map, showing only those regions where \\subref{subfig:smica_-50_50}~$|T|\\le50\\uK$ and \\subref{subfig:smica_abs_50_120}~$50\\uK\\le|T|\\le120\\uK$. The small Galactic mask used in the analysis appears as a thin horizontal gray line in the center of the maps, and the masked pixels are not included in any of the temperature ranges. In panel~\\subref{subfig:smica_abs_50_120}, the black curves mark the location of the ecliptic plane.}\n\t\\label{fig:smica_regions}\n\\end{figure}\nWe see that there is no apparent tendency for the contributions from either of these temperature ranges to be localized in specific regions of the sky. We do, however, pay special attention to the region of the ecliptic plane. As is apparent from fig.~3 of~\\cite{Planck15Maps}, the SMICA map is susceptible to contamination from foreground residuals in the region of the ecliptic. We therefore include in Fig.~\\ref{subfig:smica_abs_50_120} curves showing the location of the ecliptic plane and suggest that the number of cold spots in the ecliptic band might be unexpected. As it is not the focus of this work, we have not performed any quantitative analysis regarding the spatial distribution of hot or cold regions of the map. The maps in Fig.~\\ref{fig:smica_regions} provide\nan additional general indication that the SMICA temperature map is not anomalous. Nevertheless, we again emphasize that small foreground residuals, like those suspected to lie in the area of the ecliptic plane, while insignificant to the analysis of temperature fluctuations, can become extremely important when analyzing the CMB polarization pattern, specifically $B$-mode polarization. \n\n\n\n\\subsection{Interchangeability of $P$ and $Q$}\n\\label{sub:interchangeability_of_P_and_Q}\n\nAs has been noted above, the KL divergence is not symmetric with respect to the interchange of the distributions $P$ and $Q$ except in the limit $P \\to Q$.\nThis is a reminder of the fact that the KL divergence is not a true metric of the distance between $P$ and $Q$. So far, we have followed the common practice of taking $P$ to be the distribution of the data and $Q$ the expected distribution~\\cite{kl}. However, it is worth checking what happens when these roles are reversed. We have performed two tests involving interchange of the two distributions. First, we simply calculate the KL divergence $K(Q\\|P)$, where $P$ is again the SMICA histogram and $Q$ is the mean histogram of the FFP7 ensemble.\\footnote{Note that with the roles reversed, bins with $P_i < 5$ are now ignored and bins with small values of $Q_i$ are all counted.} This value is then compared to the ensemble of values computed with $P$ replaced by each of the FFP7 maps. The\nresulting histogram, after normalization using Eq.~\\eqref{eq:alpha_definition}, is presented in Fig.~\\ref{subfig:PQ_and_QP} together with the histogram of $K(P\\|Q)$ (presented above) as reference.\n\\begin{figure}\n\t\\centering\n\t\\subfigure[]{\n\t\t\\includegraphics[width=\\halfW]{KL_QP_hist.ps}\n\t\t\\label{subfig:PQ_and_QP}\n\t}\n\t\\subfigure[]{\n\t\t\\includegraphics[width=\\halfW]{Delta_KL_hist.ps}\n\t\t\\label{subfig:PQ_-_QP}\n\t}\n\t\\caption{\\subref{subfig:PQ_and_QP}~Histograms for the usual KL divergence $K(P\\|Q)$ (\\emph{black}) and for the reversed divergence $K(Q\\|P)$\n\t(\\emph{blue}). The red and blue vertical lines are the values for the SMICA map for the normal and reversed tests, respectively. All $K$ values have been normalized using Eq.~\\eqref{eq:alpha_definition}.\n\\subref{subfig:PQ_-_QP}\n\t~Histogram of the normalized difference $\\alpha(P\\|Q)-\\alpha(Q\\|P)$. The vertical line is the value for the SMICA map.}\n\t\\label{fig:PQQP}\n\\end{figure}\nWe can see that the two histograms and the corresponding $p$-values are similar. The value $p = 6.5\\%$ was obtained for the reversed test; the\n$p$-value for the normal test is $5.6\\%$ as stated above. It is apparent that, although similar, the reversed histogram is slightly but consistently shifted towards smaller $\\alpha$ values than the normal histogram. The value for the SMICA map is also lower for the reversed test. However, the SMICA is a single map, and the shift between histograms only indicates a statistical shift for the whole ensemble. Therefore, the second test is to examine the difference $\\Delta\\alpha = \\alpha(P\\|Q)-\\alpha(Q\\|P)$ between the normal and reversed normalized KL divergences of the \\emph{same} map. Fig.~\\ref{subfig:PQ_-_QP} shows the resulting histogram for the FFP7 ensemble together with the value for SMICA. The SMICA map shows a highly standard $\\Delta\\alpha$, yielding a $p$-value of $49.4\\%$. A similar test performed versus the Gaussian ensemble gives very similar results.\n\nThe tendency of the KL divergences to become smaller when $P$ and $Q$ are interchanged can be understood easily. Since $P$ is calculated from a single map, it tends to fluctuate more than $Q$, which is the mean of all ensemble distributions. Therefore, when $P$ appears only inside the logarithm, as is the case in the reversed $K(Q\\|P)$, the fluctuations are suppressed relative to the normal $K(P\\|Q)$. We see here that the SMICA map not only shows the expected qualitative behavior upon interchanging $P$ and $Q$, it is also quantitatively shifted by the expected amount. While the KL divergence in general is not symmetric under the interchange of $P$ and $Q$, we conclude that when testing the one-dimensional temperature distribution of the CMB on large scales, reversing the two makes little difference.\n\n\\section{Discussion}\n\\label{sec:discussion}\n\nWe have discussed the applicability of the Kullback--Leibler divergence for the assessment of departures form Gaussianity of CMB temperature\nmaps on large scales. We have illustrated this on the SMICA map, comparing it to both the set of FFP7 simulations and a set of $10^3$ Gaussian draws.\nWe have shown that it is consistent with each of these reference sets to a level of about 6\\%. We have used the KL divergence to compare the FFP7 and Gaussian reference sets and have shown that they are in good agreement. This suggests that the additional instrumental effects and foreground residuals included in the FFP7 simulations are unimportant on the scales considered here. Since the KL divergence is not symmetric in $P$ and $Q$, we have performed tests to demonstrate that their interchange has little effect on these conclusions. Finally, we have repeated these calculations using the Kolmogorov--Smirnov test. The resulting\n$p$-value of about 3\\% suggests that the differences between the two tests are not large. We note that there is no guarantee that these tests will always\ngive similar results. For example, the KL divergence is likely to be far more sensitive than the KS test for situations where there are large relative differences in the small amplitude tails of the distributions. We have also repeated all the tests on the CMB data of the 2015 release from Planck, which recently became publicly available.\\footnote{See the Planck Legacy Archive \\url{http:\/\/pla.esac.esa.int\/pla\/}.} The results on the 2015 data set are in very good agreement with those reported here.\n\nThe difficulty in devising tests for the assessment of non-Gaussianity of the temperature and polarization maps of the CMB lies in our ignorance of\nthe nature of the non-Gaussian residuals from foregrounds and systematic effects that could propagate to the maps. In such circumstances, it seems advisable to adopt a procedure that uses as much information in the maps as possible. With its connection to the intrinsic and cross-entropy of the distributions $P$ and $Q$,\nthe KL divergence would appear to be the natural choice. Given the correlations between the pixels of the CMB, a consequence of a random draw in harmonic space, this is not necessarily the case. However, the non-locality of the KL divergence and its sensitivity to the tails of the distributions still suggest that it is a valuable complement to the KS test and might be a useful alternative. Indeed, one should utilize a variety of methods and tests to identify possible contamination of the cosmological product. Obviously, \\emph{any\\\/} suggestion of an anomalous result would indicate the need for more sophisticated analyses to assess the quality of the CMB maps.\n\n\n\\acknowledgments\n\nWe would like to thank P. Naselsky for his enthusiastic support and encouragement of the work reported in this paper and for many valuable scientific discussions.\n\nThis work is based on observations obtained with Planck,\\footnote{\\url{http:\/\/www.esa.int\/Planck}} an ESA\nscience mission with instruments and contributions directly funded by ESA\nMember States, NASA, and Canada. The development of Planck has been\nsupported by: ESA; CNES and CNRS \/ INSU-IN2P3-INP (France); ASI, CNR, and\nINAF (Italy); NASA and DoE (USA); STFC and UKSA (UK); CSIC, MICINN and JA\n(Spain); Tekes, AoF and CSC (Finland); DLR and MPG (Germany); CSA (Canada);\nDTU Space (Denmark); SER\/SSO (Switzerland); RCN (Norway); SFI (Ireland);\nFCT\/MCTES (Portugal); and PRACE (EU). A description of the Planck Collaboration and a list of its members,\nincluding the technical or scientific activities in which they have been\ninvolved, can be found at the Planck web page.\\footnote{\\url{http:\/\/www.cosmos.esa.int\/web\/planck\/planck-collaboration}}\n\nWe acknowledge the use of the NASA Legacy Archive for Microwave Background Data Analysis (LAMBDA). Our data analysis made use of the GLESP package\\footnote{\\url{http:\/\/www.glesp.nbi.dk\/}}~\\cite{Glesp}, and of \\textsc{HEALPix}~\\cite{HEALPix}. This work is supported in part by Danmarks Grundforskningsfond which allowed the establishment of the Danish Discovery Center, FNU grants 272-06-0417, 272-07-0528 and 21-04-0355, the National Natural Science Foundation of China (Grant No.\\ 11033003), the National Natural Science Foundation for Young Scientists of China (Grant No.\\ 11203024) and the Youth Innovation Promotion Association, CAS.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nAs part of the proof of his eponymous theorem~\\cite{Szemeredi75} on arithmetic progressions in dense sets of integers, Szemer\\'edi developed (a variant of what is now known as) the graph {\\em regularity lemma}~\\cite{Szemeredi78}.\nThe lemma roughly states that the vertex set of every graph can be partitioned into a bounded number of parts such that\nalmost all the bipartite graphs induced by pairs of parts in the partition are quasi-random.\nIn the past four decades this lemma has become one of the (if not the) most powerful tools in extremal combinatorics, with applications in many other areas of mathematics. We refer the reader to~\\cite{KomlosShSiSz02,RodlSc10}\nfor more background on the graph regularity lemma, its many variants and its numerous applications.\n\nPerhaps the most important and well-known application of the graph regularity lemma is the original\nproof of the {\\em triangle removal lemma}, which states that if an $n$-vertex graph $G$ contains only $o(n^3)$ triangles, then one can turn $G$ into a triangle-free graph by removing only $o(n^2)$ edges (see \\cite{ConlonFox13} for more details).\nIt was famously observed by Ruzsa and Szemer\\'edi~\\cite{RuzsaSz76} that the triangle removal lemma implies Roth's theorem~\\cite{Roth54}, the special case of Szemer\\'edi's theorem for $3$-term arithmetic progressions.\nThe problem of extending the triangle removal lemma to the hypergraph\nsetting was raised by Erd\\H{o}s, Frankl and R\\\"odl~\\cite{ErdosFrRo86}. One of the main motivations for obtaining such a result was the observation of Frankl and R\\\"odl~\\cite{FrankRo02} (see also~\\cite{Solymosi04}) that such a result would allow one to extend the Ruzsa--Szemer\\'edi~\\cite{RuzsaSz76} argument and thus obtain an alternative proof of Szemer\\'edi's theorem for progressions of arbitrary length.\n\nThe quest for a hypergraph regularity lemma, which would allow one to prove a hypergraph removal lemma, took about 20 years.\nThe first milestone was the result of Frankl and R\\\"odl~\\cite{FrankRo02}, who obtained a regularity lemma for $3$-uniform hypergraphs. About 10 years later, the approach of~\\cite{FrankRo02} was extended to hypergraphs of arbitrary uniformity by R\\\"odl, Skokan, Nagle and Schacht~\\cite{NagleRoSc06, RodlSk04}.\nAt the same time, Gowers~\\cite{Gowers07} obtained an alternative version of the regularity lemma for $k$-uniform hypergraphs (from now on we will use $k$-graphs instead of $k$-uniform hypergraphs).\nShortly after, Tao~\\cite{Tao06} and R\\\"odl and Schacht~\\cite{RodlSc07,RodlSc07-B} obtained two more versions of the lemma.\n\nAs it turned out, the main difficulty with obtaining a regularity lemma for $k$-graphs was defining the correct notion of hypergraph regularity that would:\n$(i)$ be strong enough to allow one to prove a counting lemma, and\n$(ii)$ be weak enough to be satisfied by every hypergraph (see the discussion in~\\cite{Gowers06} for more on this issue).\nAnd indeed, the above-mentioned variants of the hypergraph regularity lemma rely on four different notions of quasi-randomness, which\nto this date are still not known to be equivalent\\footnote{This should be contrasted with the setting of graphs in which\n(almost) all notions of quasi-randomness are not only known to be equivalent but even effectively equivalent. See e.g.~\\cite{ChungGrWi89}.} (see~\\cite{NaglePoRoSc09} for some partial results). What all of these proofs {\\em do} have in common however,\nis that they supply only Ackermann-type bounds for the size of a regular partition.\\footnote{Another variant of the hypergraph regularity lemma was obtained in~\\cite{ElekSz12}. This approach does not supply any quantitative bounds.} \nMore precisely, if we let $\\Ack_1(x)=2^x$ and then define $\\Ack_k(x)$ to be the $x$-times iterated\\footnote{$\\Ack_2(x)$ is thus a tower of exponents of height $x$,\n\t$\\Ack_3(x)$ is the so-called wowzer function, etc.} version of $\\Ack_{k-1}$, then all the above proofs guarantee to produce a regular partition of a $k$-graph whose order can be bounded from above by an $\\Ack_k$-type function.\n\nOne of the most important applications of the $k$-graph regularity lemma was that it gave the first explicit\nbounds for the multidimensional generalization of Szemer\\'edi's theorem, see~\\cite{Gowers07}. The original proof of this result, obtained by Furstenberg and Katznelson~\\cite{FurstenbergKa78}, relied on Ergodic Theory and thus supplied no quantitative bounds at all.\nExamining the reduction between these theorems~\\cite{Solymosi04} reveals that if one could improve the Ackermann-type bounds for the $k$-graph regularity\nlemma, by obtaining (say) $\\Ack_{k_0}$-type upper bounds (for all $k$), then one would obtain the first primitive recursive bounds for the\nmultidimensional generalization of Szemer\\'edi's theorem. Let us note that obtaining such bounds just for\nvan der Waerden's theorem~\\cite{Shelah89} and Szemer\\'edi's theorem~\\cite{Szemeredi75} (which are two special case) were open problems for many decades till\nthey were finally solved by Shelah~\\cite{Shelah89} and Gowers~\\cite{Gowers01}, respectively.\nFurther applications of the $k$-graph regularity lemma (and the hypergraph removal lemma in particular) are described in~\\cite{RodlNaSkScKo05} and~\\cite{RodlTeScTo06} as well as in R\\\"odl's recent ICM survey~\\cite{Rodl14}.\n\n\nA famous result of Gowers~\\cite{Gowers97} states that the $\\Ack_2$-type upper bounds for graph regularity\nare unavoidable. Several improvements~\\cite{FoxLo17},\nvariants~\\cite{ConlonFo12,KaliSh13,MoshkovitzSh18} and simplifications~\\cite{MoshkovitzSh16} of Gowers' lower bound were recently obtained, but no analogous lower bound was derived even for $3$-graph regularity.\nThe numerous applications of the hypergraph regularity lemma naturally lead to the question of whether one can improve upon the Ackermann-type\nbounds mentioned above and obtain primitive recursive bounds\nfor the $k$-graph regularity lemma. Tao~\\cite{Tao06-h} predicted that the answer to this question is negative, in the sense that one cannot obtain better than $\\Ack_k$-type upper bounds for the $k$-graph regularity lemma for every $k \\ge 2$.\nThe main result presented here and in the followup \\cite{MSk} confirms this prediction.\n\n\n\\begin{theo}{\\bf[Main result, informal statement]}\\label{thm:main-informal}\nThe following holds for every $k\\geq 2$: every regularity lemma for $k$-graphs satisfying some\nmild conditions can only guarantee to produce partitions of size bounded by an $\\Ack_k$-type function.\n\\end{theo}\n\nIn this paper we will focus on proving the key ingredient needed for obtaining Theorem~\\ref{thm:main-informal},\nstated as Lemma~\\ref{theo:core} in Subsection~\\ref{subsec:overview}, and on showing how it can be used in order to prove Theorem~\\ref{thm:main-informal}\nfor $k=3$. In a nutshell, the key idea is to use the graph construction given by Lemma~\\ref{theo:core} in order\nto construct a $3$-graph by taking a certain ``product'' of two graphs that are hard for graph regularity, in order to get\na $3$-graph that is hard for $3$-graph regularity. See the discussion following Lemma~\\ref{theo:core} in Subsection~\\ref{subsec:overview}.\nDealing with $k=3$ in this paper will allow us to present all the new ideas needed in order to actually prove Theorem~\\ref{thm:main-informal}\nfor arbitrary $k$, in the slightly friendlier setting of $3$-graphs.\nIn a followup paper~\\cite{MSk}, we will show how Lemma~\\ref{theo:core} can be used in order to prove\nTheorem~\\ref{thm:main-informal} for all $k \\ge 2$.\n\nIn this paper we will also show how to derive from Theorem~\\ref{thm:main-informal} tight lower bounds for\nthe $3$-graph regularity lemmas due to Frankl and R\\\"odl~\\cite{FrankRo02} and to Gowers~\\cite{Gowers06}.\n\n\\begin{coro}\\label{coro:FR-LB}\nThere is an $\\Ack_3$-type lower bound for the $3$-graph regularity lemmas of Frankl and R\\\"odl~\\cite{FrankRo02} and of Gowers~\\cite{Gowers06}.\n\\end{coro}\n\nIn \\cite{MSk} we will show how to derive from Theorem~\\ref{thm:main-informal} a tight lower bound for\nthe $k$-graph regularity lemma due to R\\\"odl and Schacht~\\cite{RodlSc07}.\n\n\\begin{coro}\\label{coro:RS-LB}\nThere is an $\\Ack_k$-type lower bound for the $k$-graph regularity lemma of R\\\"odl and Schacht~\\cite{RodlSc07}.\n\\end{coro}\n\n\nBefore getting into the gory details of the proof, let us informally discuss what we think are some\ninteresting aspects of the proof of Theorem \\ref{thm:main-informal}.\n\n\\paragraph{Why is it hard to ``step up''?}\nThe reason why the upper bound for graph regularity is of tower-type\nis that the process of constructing a regular partition of a graph proceeds by a sequence of steps, each increasing the size of the partition exponentially.\nThe main idea behind Gowers' lower bound for graph regularity~\\cite{Gowers97} is in ``reverse engineering'' the proof of the upper bound; in other words,\nin showing that (in some sense) the process of building the partition using a sequence of exponential refinements is unavoidable.\nNow, a common theme in all proofs of the hypergraph regularity lemma\nfor $k$-graphs is that they proceed by induction on $k$; that is, in the process of constructing a regular\npartition of the input $k$-graph $H$, the proof applies the $(k-1)$-graph regularity lemma on certain $(k-1)$-graphs\nderived from $H$. This is why one gets $\\Ack_k$-type upper bounds. So with~\\cite{Gowers97} in mind, one might guess\nthat in order to prove a matching lower bound one should ``reverse engineer'' the proof of the upper bound and show that such a process is unavoidable. However, this turns out to be false! As we argued in~\\cite{MoshkovitzSh18}, in order to prove an {\\em upper bound} for (say) $3$-graph regularity it is in fact enough to iterate a relaxed version of graph regularity which we call the ``sparse regular approximation lemma'' (SRAL for short).\nTherefore, in order to prove an $\\Ack_3$-type {\\em lower bound} for $3$-graph\nregularity one cannot simply ``step up'' an $\\Ack_2$-type lower bound for graph regularity. Indeed, a necessary condition\nwould be to prove an $\\Ack_2$-type lower bound for SRAL. See also the discussion following Lemma~\\ref{theo:core} in Subsection \\ref{subsec:overview}\non how do we actually use a graph construction in order to get a $3$-graph construction.\n\n\n\\paragraph{A new notion of graph\/hypergraph regularity:}\nIn a recent paper \\cite{MoshkovitzSh18} we proved an $\\Ack_2$-type\nlower bound for SRAL.\nAs it turned out, even this lower bound was not enough to allow us to step up the graph lower bound\ninto a $3$-graph lower bound. To remedy this, in the present paper we introduce an even weaker notion of graph\/hypergraph regularity\nwhich we call $\\langle \\d \\rangle$-regularity. This notion seems to be right at the correct level of ``strength'';\non the one hand, it is strong enough to allow one to prove $\\Ack_{k-1}$-type lower bounds for $(k-1)$-graph regularity, while at the\nsame time weak enough to allow one to induct, that is, to use it in order to then prove $\\Ack_{k}$-type lower bounds for $k$-graph regularity.\nAnother critical feature of our new notion of hypergraph regularity is that it has (almost) nothing to do with hypergraphs!\nA disconcerting aspect of all proofs of the hypergraph regularity lemma is that they involve a very complicated nested\/inductive structure.\nFurthermore, one has to introduce an elaborate hierarchy of constants that controls how regular one level of the partition is compared to\nthe previous one. What is thus nice about our new notion is that it involves only various kinds of instances of graph $\\langle \\d \\rangle$-regularity.\nAs a result, our proof is (relatively!) simple.\n\n\\paragraph{How do we find witnesses for $3$-graph irregularity?}\nThe key idea in Gowers' lower bound~\\cite{Gowers97} for graph regularity was in constructing a graph $G$, based on a sequence of\npartitions ${\\cal P}_1,{\\cal P}_2,\\ldots$ of $V(G)$, with the following\ninductive property: if a vertex partition $\\Z$ refines ${\\cal P}_i$\nbut does not refine ${\\cal P}_{i+1}$ then $\\Z$ is not $\\epsilon$-regular.\nThe key step of the proof of~\\cite{Gowers97} is in finding witnesses showing\nthat pairs of clusters of $\\Z$ are irregular. The main difficulty in extending this\nstrategy to $k$-graphs already reveals itself in the setting of $3$-graphs. In a nutshell, while\nin graphs, a witness to irregularity of a pair of clusters $A,B \\in \\Z$ is\n{\\em any} pair of large subsets $A' \\sub A$ and $B' \\sub B$, in the setting of $3$-graphs\nwe have to find three large edge-sets (called a {\\em triad}, see Section~\\ref{sec:FR}) that\nhave an additional property: they must together form a graph containing many triangles.\nIt thus seems quite hard to extend Gowers' approach already to the setting of $3$-graphs. By working\nwith the much weaker notion of $\\langle \\d \\rangle$-regularity, we circumvent this\nissue since two of the edges sets in our version of a triad are always complete bipartite graphs.\nSee Subsection~\\ref{subsec:definitions}.\n\n\\paragraph{What is then the meaning of Theorem \\ref{thm:main-informal}?}\nOur main result, stated formally as Theorem~\\ref{theo:main}, establishes an $\\Ack_3$-type lower bound\nfor $\\langle \\d \\rangle$-regularity of $3$-graphs, that is, for a specific new\nversion of the hypergraph regularity lemma. Therefore, we immediately get $\\Ack_3$-type lower bounds for\nany $3$-graph regularity lemma which is at least as strong as our new lemma, that is, for any lemma\nwhose requirements\/guarantees imply those that are needed in order to satisfy our new notion of regularity.\nIn particular, we will prove Corollary \\ref{coro:FR-LB} by showing that the regularity notions\nused in these lemmas are at least as strong as $\\langle \\d \\rangle$-regularity.\n\nIn \\cite{MSk} we will prove Theorem~\\ref{thm:main-informal} in its full generality by extending Theorem~\\ref{theo:main}\nto arbitrary $k$-graphs. This proof, though technically more involved, will be\nquite similar at its core to the way we derive Theorem~\\ref{theo:main} from Lemma~\\ref{theo:core} in the present paper.\nThe deduction of Corollary \\ref{coro:RS-LB}, which appears in \\cite{MSk}, will also turn out to be quite similar to the way\nCorollary \\ref{coro:FR-LB} is derived from Theorem~\\ref{theo:main} in the present paper\n\n\\paragraph{How strong is our lower bound?} Since Theorem \\ref{thm:main-informal} gives a lower bound for\n$\\langle \\d \\rangle$-regularity and Corollaries \\ref{coro:FR-LB} and \\ref{coro:RS-LB} show that\nthis notion is at least as weak as previously used notions of regularity, it is natural to ask:\n$(i)$ is this notion equivalent to one of the other notions? $(ii)$ is this notion strong enough for proving the\nhypergraph removal lemma, which was one of the main reasons for developing\nthe hypergraph regularity lemma? We will prove that the answer to both questions is {\\em negative} by showing that already for\ngraphs, $\\langle \\d \\rangle$-regularity (for $\\d$ a fixed constant) is not strong enough even for proving the triangle removal lemma.\nThis of course makes our lower bound even stronger as it already applies to a very weak notion of regularity.\nIn a nutshell, the proof proceeds by first taking a random tripartite graph, showing (using routine probabilistic\narguments) that with high probability the graph is $\\langle \\d \\rangle$-regular yet contains a small number of triangles.\nOne then shows that removing these triangles, and then taking a blowup of the resulting graph, gives a triangle-free graph of positive density that is $\\langle \\d \\rangle$-regular. The full details will appear in \\cite{MSk}.\n\n\n\\paragraph{How tight is our bound?} Roughly speaking, we will show that for a $k$-graph with $pn^k$ edges,\nevery $\\langle \\d \\rangle$-regular partition has order at least $\\Ack_k(\\log 1\/p)$. In a recent paper \\cite{MoshkovitzSh16}\nwe proved that in graphs, one can prove a matching $\\Ack_2(\\log 1\/p)$ upper bound, even for a slightly stronger notion than $\\langle \\d \\rangle$-regularity.\nThis allowed us to obtain a new proof of Fox's $\\Ack_2(\\log 1\/\\epsilon)$ upper bound for the graph removal lemma \\cite{Fox11} (since the stronger notion allows to count small subgraphs). We believe that it should\nbe possible to match our lower bounds with $\\Ack_k(\\log 1\/p)$ upper bounds (even for a slightly stronger notion analogous to the one used in \\cite{MoshkovitzSh16}). We think that it should be possible to deduce from such an upper bound an $\\Ack_k(\\log 1\/\\epsilon)$ upper bound\nfor the $k$-graph removal lemma. The best known bounds for this problem are (at least) $\\Ack_k(\\poly(1\/\\epsilon))$.\n\n\n\n\n\\subsection{Paper overview}\n\nIn Section~\\ref{sec:define} we will first define the new notion of hypergraph regularity, which we term $\\langle \\d \\rangle$-regularity,\nfor which we will prove our main lower bound. We will then give the formal version of Theorem~\\ref{thm:main-informal} (see Theorem~\\ref{theo:main}).\nThis will be followed by the statement of our core technical result, Lemma~\\ref{theo:core},\nand an overview of how this technical result is used in the proof of Theorem~\\ref{theo:main}.\nThe proof of Theorem~\\ref{theo:main} appears in Section~\\ref{sec:LB}. \nWe refer the reader to~\\cite{MSk} for the proof of Lemma~\\ref{theo:core}.\nIn Section \\ref{sec:FR} we prove Corollary~\\ref{coro:FR-LB}.\nIn Appendix~\\ref{sec:FR-appendix} we give the proof of certain technical claims missing from Section~\\ref{sec:FR}.\n\n\n\n\n\n\n\\section{$\\langle \\d \\rangle$-regularity and Proof Overview}\\label{sec:define}\n\n\nFormally, a \\emph{$3$-graph} is a pair $H=(V,E)$, where $V=V(H)$ is the vertex set and $E=E(H) \\sub \\binom{V}{3}$ is the edge set of $H$.\nThe number of edges of $H$ is denoted $e(H)$ (i.e., $e(H)=|E|$).\nThe $3$-graph $H$ is \\emph{$3$-partite} on (disjoint) vertex classes $(V_1,V_2,V_3)$ if every edge of $H$ has a vertex from each $V_i$.\nThe \\emph{density} of a $3$-partite $3$-graph $H$ is $e(H)\/\\prod_{i=1}^3 |V_i|$.\nFor a bipartite graph $G$, the set of edges of $G$ between disjoint vertex subsets $A$ and $B$ is denoted by $E_G(A,B)$; the density of $G$ between $A$ and $B$ is denoted by $d_G(A,B)=e_G(A,B)\/|A||B|$, where $e_G(A,B)=|E_G(A,B)|$. We use $d(A,B)$ if $G$ is clear from context.\nWhen it is clear from context, we sometimes identify a hypergraph with its edge set. In particular, we will write $V_1 \\times V_2$ for the complete bipartite graph on vertex classes $(V_1,V_2)$.\nFor partitions $\\P,\\Q$ of the same underlying set, we say that $\\Q$ \\emph{refines} $\\P$, denoted $\\Q \\prec \\P$, if every member of $\\Q$ is contained in a member of $\\P$.\nWe say that $\\P$ is \\emph{equitable} if all its members have the same size.\\footnote{In a regularity lemma one allows the parts to differ in size by at most $1$ so that it applies to all (hyper-)graphs. For our lower bound this is unnecessary.}\nWe use the notation $x \\pm \\e$ for a number lying in the interval $[x-\\e,\\,x+\\e]$.\n\nIn the following definition, and in the rest of the paper, we will sometimes identify a graph or a $3$-graph with its edge set when the vertex set is clear from context.\n\n\\begin{definition}[$2$-partition]\\label{def:2-partition}\nA \\emph{$2$-partition} $(\\Z,\\E)$ on a vertex set $V$ consists of a partition $\\Z$ of $V$ and a family of edge disjoint bipartite graphs\n$\\E$ so that:\n\\begin{itemize}\n\\item Every $E \\in \\E$ is a bipartite graph\nwhose two vertex sets are distinct $Z,Z' \\in \\Z$.\n\\item For every $Z \\neq Z' \\in \\Z$, the complete bipartite graph $Z \\times Z'$ is the union of graphs from $\\E$.\n\\end{itemize}\n\n\n\n\n\n\\end{definition}\n\nPut differently, a $2$-partition consists of vertex partition $\\Z$ and a collection of bipartite graphs $\\E$ such that\n$\\E$ is a refinement of the collection of complete bipartite graphs $\\{Z \\times Z' : Z \\neq Z' \\in \\Z \\}$.\n\n\n\\subsection{$\\langle \\d \\rangle$-regularity of graphs and hypergraphs}\\label{subsec:definitions}\n\nIn this subsection we define our new\\footnote{For $k=3$, related notions of regularity were studied in~\\cite{ReiherRoSc16,Towsner17}.} notion of $\\langle\\d\\rangle$-regularity, first for graphs and then for $3$-graphs in Definition~\\ref{def:k-reg} below.\nLet us first recall Szemer\\'edi's notion of $\\epsilon$-regularity.\nA bipartite graph on $(A,B)$ is \\emph{$\\e$-regular} if for all subsets $A' \\sub A$, $B' \\sub B$ with $|A'|\\ge\\e|A|$, $|B'|\\ge\\e|B|$ we have $|d(A',B') -d(A,B)| \\le \\e$.\nA vertex partition $\\P$ of a graph is $\\e$-regular\nif\nthe bipartite graph induced on each but at most $\\e|\\P|^2$ of the pairs $(A,B)$ with $A \\neq B \\in \\P$ is $\\e$-regular.\nSzemer\\'edi's graph regularity lemma says that every graph\nhas an $\\e$-regular equipartition of order at most some $\\Ack_2(\\poly(1\/\\e))$.\nWe now introduce a weaker notion of graph regularity which we will use throughout the paper.\n\\begin{definition}[graph $\\langle\\d\\rangle$-regularity]\\label{def:star-regular}\n\tA bipartite graph $G$ on $(A,B)$ is \\emph{$\\langle \\d \\rangle$-regular} if for all subsets $A' \\sub A$, $B' \\sub B$ with $|A'| \\ge \\d|A|$, $|B'|\\ge\\d|B|$ we have $d_G(A',B') \\ge \\frac12 d_G(A,B)$.\\\\\n\tA vertex partition $\\P$ of a graph $G$ is \\emph{$\\langle \\d \\rangle$-regular}\n\n\tif one can add\/remove at most $\\d \\cdot e(G)$ edges so that the bipartite graph induced on each $(A,B)$ with $A \\neq B \\in \\P$ is $\\langle \\d \\rangle$-regular.\n\n\\end{definition}\n\nFor the reader worried that in Definition~\\ref{def:star-regular} we merely replaced the $\\e$\nfrom the definition of $\\e$-regularity with $\\d$, we refer to the discussion following Theorem~\\ref{theo:main} below.\n\n\n\n\n\nThe definition of $\\langle\\d\\rangle$-regularity for hypergraphs involves the $\\langle\\d\\rangle$-regularity notion for graphs, applied to certain auxiliary graphs which are defined as follows.\n\n\n\n\n\\begin{definition}[The auxiliary graph $G_{H}^i$]\\label{def:aux}\nFor a $3$-partite $3$-graph $H$ on vertex classes $(V_1,V_2,V_3)$,\nwe define a bipartite graph $G_{H}^1$ on the vertex classes $(V_2 \\times V_3,\\,V_1)$ by\n$$E(G_{H}^1) = \\big\\{ ((v_2,v_3),v_1) \\,\\big\\vert\\, (v_1,v_2,v_3) \\in E(H) \\big\\} \\;.$$\nThe graphs $G_{H}^2$ and $G_{H}^3$ are defined in an analogous manner.\n\\end{definition}\n\nImportantly, for a $2$-partition (as defined in Definition~\\ref{def:2-partition}) to be $\\langle\\d\\rangle$-regular it must first satisfy a requirement on the regularity of its parts.\n\n\n\\begin{definition}[$\\langle\\d\\rangle$-good partition]\\label{def:k-good}\nA $2$-partition $(\\Z,\\E)$ on $V$ is \\emph{$\\langle\\d\\rangle$-good} if all bipartite graphs in $\\E$\n(between any two distinct vertex clusters of $\\Z$) are $\\langle \\d \\rangle$-regular.\n\\end{definition}\n\n\n\n\n\n\nFor a $2$-partition $(\\Z,\\E)$ of a $3$-partite $3$-graph on vertex classes $(V_1,V_2,V_3)$ with $\\Z \\prec \\{V_1,V_2,V_3\\}$, for every $1 \\le i \\le 3$ we denote $\\Z_i = \\{Z \\in \\Z \\,\\vert\\, Z \\sub V_i\\}$, and we denote $\\E_i = \\{E \\in \\E \\,\\vert\\, E \\sub V_j \\times V_k\\}$ where $\\{i,j,k\\}=\\{1,2,3\\}$.\nSo for example, $\\E_1$ is thus a partition of $V_2 \\times V_3$.\n\n\\begin{definition}[$\\langle\\d\\rangle$-regular partition]\\label{def:k-reg}\nLet $H$ be a $3$-partite $3$-graph on vertex classes $(V_1,V_2,V_3)$\nand $(\\Z,\\E)$ be a $\\langle \\d \\rangle$-good $2$-partition with $\\Z \\prec \\{V_1,V_2,V_3\\}$.\nWe say that $(\\Z,\\E)$ is a \\emph{$\\langle \\d \\rangle$-regular} partition of $H$ if\nfor every $1 \\le i \\le 3$,\n$\\E_i \\cup \\Z_i$ is a $\\langle \\d \\rangle$-regular partition of $G_H^i$.\n\n\n\n\n\n\\end{definition}\n\n\n\n\n\n\n\n\n\n\t\t\n\t\t\n\n\n\\subsection{Formal statement of the main result}\n\nWe are now ready to formally state our tight lower bound for $3$-graph $\\langle \\d \\rangle$-regularity (the formal version of\nTheorem~\\ref{thm:main-informal} above for $k=3$). Recall that we define the {\\em tower} functions $\\twr(x)$ to be a tower of exponents of height $x$,\nand then define the {\\em wowzer} function $\\wow(x)$ to be the $x$-times iterated tower function, that is\n$\\wow(x)= \\underbrace{\\twr(\\twr(\\cdots(\\twr(1))\\cdots))}_{x \\text{ times}}$.\n\n\\begin{theo}[Main result]\\label{theo:main}\nFor every $s \\in \\N$ there is a $3$-partite $3$-graph $H$ on vertex classes of equal size and of density at least $2^{-s}$,\nand a partition $\\V_0$ of $V(H)$ with $|\\V_0| \\le 2^{300}$, such that if\n$(\\Z,\\E)$ is a $\\langle 2^{-73} \\rangle$-regular partition of $H$ with\n$\\Z \\prec \\V_0$ then $|\\Z| \\ge \\wow(s)$.\n\n\n\\end{theo}\n\nLet us draw the reader's attention to an important and perhaps surprising aspect of Theorem~\\ref{theo:main}.\nAll the known tower-type lower bounds for graph regularity depend on the error parameter $\\epsilon$,\nthat is, they show the existence of graphs $G$ with the property that every $\\epsilon$-regular partition of $G$ is of order at least $\\Ack_2(\\poly(1\/\\e))$.\nThis should be contrasted with the fact that our lower bounds for $\\langle \\d \\rangle$-regularity holds for a {\\em fixed} error parameter $\\delta$.\nIndeed, instead of the dependence on the error parameter, our lower bound depends on the {\\em density} of the graph.\nThis delicate difference makes it possible for us to prove Theorem~\\ref{theo:main} by iterating the construction described in the next subsection.\n\n\n\\subsection{The core construction and proof overview}\\label{subsec:overview}\n\nThe graph construction in Lemma~\\ref{theo:core} below is the main technical result we will need in order to prove Theorem~\\ref{theo:main}.\nWe will first need to define ``approximate'' refinement (a notion that goes back to Gowers~\\cite{Gowers97}).\n\\begin{definition}[Approximate refinements]\nFor sets $S,T$ we write $S \\sub_\\b T$ if $|S \\sm T| < \\b|S|$.\nFor a partition $\\P$ we write $S \\in_\\b \\P$ if $S \\sub_\\b P$ for some $P \\in \\P$.\nFor partitions $\\P,\\Q$ of the same set of size $n$ we write $\\Q \\prec_\\b \\P$ if\n$$\\sum_{\\substack{Q \\in \\Q\\colon\\\\Q \\notin_\\b \\P}} |Q| \\le \\b n \\;.$$\n\\end{definition}\nNote that for $\\Q$ equitable, $\\Q \\prec_\\b \\P$ if and only if\nall but at most $\\b|\\Q|$ parts $Q \\in \\Q$ satisfy $Q \\in_\\b \\P$.\nWe note that throughout the paper we will only use approximate refinements with $\\b \\le 1\/2$, and so if $S \\in_\\b \\P$ then $S \\sub_\\b P$ for a unique $P \\in \\P$.\n\nWe stress that in Lemma~\\ref{theo:core} below we only use notions related to graphs. In particular, $\\langle \\d \\rangle$-regularity refers to Definition~\\ref{def:star-regular}.\n\n\\begin{lemma}[Core construction]\\label{theo:core}\nLet $\\Lside$ and $\\Rside$ be disjoint sets. Let\n$\\L_1 \\succ \\cdots \\succ \\L_s$ and $\\R_1 \\succ \\cdots \\succ \\R_s$ be two sequences of $s$ successively refined equipartitions of $\\Lside$ and $\\Rside$, respectively,\nthat satisfy for every $i \\ge 1$ that:\n\\begin{enumerate}\n\\item\\label{item:core-minR}\n$|\\R_i|$ is a power of $2$ and $|\\R_1| \\ge 2^{200}$,\n\\item\\label{item:core-expR} $|\\R_{i+1}| \\ge 4|\\R_i|$ if $i < s$,\n\\item\\label{item:core-expL} $|\\L_i| = 2^{|\\R_i|\/2^{i+10}}$.\n\\end{enumerate}\nThen there exists a sequence of $s$ successively refined edge equipartitions $\\G_1 \\succ \\cdots \\succ \\G_s$ of $\\Lside \\times \\Rside$ such that for every $1 \\le j \\le s$, $|\\G_j|=2^j$, and the following holds for every $G \\in \\G_j$ and $\\d \\le 2^{-20}$.\nFor every $\\langle \\d \\rangle$-regular partition $\\P \\cup \\Q$ of $G$, where $\\P$ and $\\Q$ are partitions of $\\Lside$ and $\\Rside$, respectively, and every $1 \\le i \\le j$, if $\\Q \\prec_{2^{-9}} \\R_{i}$ then $\\P \\prec_{\\g} \\L_{i}$ with $\\g = \\max\\{2^{5}\\sqrt{\\d},\\, 32\/\\sqrt[6]{|\\R_1|} \\}$.\n\\end{lemma}\n\\begin{remark}\\label{remequi}\nEvery $G \\in \\G_j$ is a bipartite graph of density $2^{-j}$ since $\\G_j$ is equitable.\n\\end{remark}\n\nAs mentioned before, the proof of Lemma~\\ref{theo:core} appears in~\\cite{MSk}.\nLet us end this section by explaining the role Lemma~\\ref{theo:core} plays in the proof of Theorem \\ref{theo:main}.\n\n\\paragraph{Using graphs to construct $3$-graphs:}\nPerhaps the most surprising aspect of the proof of Theorem~\\ref{theo:main} is that in order to construct a $3$-graph we also use the graph construction of Lemma~\\ref{theo:core} in a somewhat unexpected way.\nIn this case, $\\Lside$ will be a complete bipartite graph and the $\\L_i$'s will be partitions of this complete bipartite graph themselves given by another application of Lemma~\\ref{theo:core}.\nThe partitions will be of wowzer-type growth, and the second application of Lemma~\\ref{theo:core} will ``multiply'' the graph partitions (given by the $\\L_i$'s) to\ngive a partition of the complete $3$-partite $3$-graph into $3$-graphs that are hard for $\\langle \\d \\rangle$-regularity.\nWe will take $H$ in Theorem~\\ref{theo:main} to be an arbitrary $3$-graph in this partition.\n\n\n\n\n\n\\paragraph{Why is Lemma~\\ref{theo:core} one-sided?}\nAs is evident from the statement of Lemma~\\ref{theo:core}, it is one-sided in nature; that is, under the premise that the partition $\\Q$\nrefines $\\R_i$ we may conclude that $\\P$ refines $\\L_i$.\nIt is natural to ask if one can do away with this assumption, that is, be able to show that\nunder the same assumptions $\\Q$ refines $\\R_i$ and $\\P$ refines $\\L_i$.\nAs we mentioned in the previous item, in order to prove a wowzer-type lower bound for $3$-graph\nregularity we have to apply Lemma~\\ref{theo:core} with a sequence of partitions that grows as a wowzer-type function.\nNow, in this setting, Lemma~\\ref{theo:core} does not hold without the one-sided assumption, because if it did, then\none would have been able to prove a wowzer-type lower bound for graph $\\langle \\d \\rangle$-regularity, and hence also for Szemer\\'edi's regularity lemma.\nPut differently, if one wishes to have a construction that holds with arbitrarily fast growing partition sizes, then\none has to introduce the one-sided assumption.\n\n\\paragraph{How do we remove the one-sided assumption?}\nThe proof of Theorem \\ref{theo:main} proceeds by first proving a one-sided version of Theorem \\ref{theo:main},\nstated as Lemma~\\ref{lemma:ind-k}. In order to get a construction that does not require such a one-sided assumption,\nwe will need one final trick; we will take $6$ clusters of\nvertices and arrange $6$ copies\nof this one-sided construction\nalong the $3$-edges of a cycle. This will give us a ``circle of implications'' that will eliminate the one-sided assumption.\nSee Subsection \\ref{subsec:pasting}.\n\n\n\n\n\n\n\n\\section{Proof of Theorem~\\ref{theo:main}}\\label{sec:LB}\n\n\n\n\\renewcommand{\\k}{r}\n\n\\renewcommand{\\t}{t}\n\\newcommand{\\w}{w}\n\n\\newcommand{\\GG}{\\mathbf{G}}\n\\newcommand{\\FF}{\\mathbf{F}}\n\\newcommand{\\VV}{\\mathbf{V}}\n\n\n\\renewcommand{\\Hy}[1]{H_{{#1}}}\n\\renewcommand{\\A}{A}\n\n\\newcommand{\\subs}{\\subset_*}\n\\newcommand{\\pad}{P}\n\\renewcommand{\\K}{\\mathcal{K}}\n\\newcommand{\\U}{U}\n\n\n\\renewcommand{\\k}{k}\n\n\\renewcommand{\\K}{\\mathcal{K}}\n\n\\renewcommand{\\r}{k}\n\n\nThe purpose of this section is to prove the main result, Theorem~\\ref{theo:main}.\nThis section is self-contained save for the application of Lemma~\\ref{theo:core}.\nThe key step of the proof, stated as Lemma~\\ref{lemma:ind-k} and proved in Subsection \\ref{subsec:key}, relies on a subtle construction that uses Lemma~\\ref{theo:core} twice. This lemma only gives a ``one-sided'' lower bound for $3$-graph regularity, in the spirit of Lemma~\\ref{theo:core}.\nIn Subsection~\\ref{subsec:pasting} we show how to use Lemma~\\ref{lemma:ind-k} in order to complete the proof of Theorem~\\ref{theo:main}.\n\nWe first observe a simple yet crucial property of $2$-partitions, stated as Claim~\\ref{claim:uniform-refinement} below, which we will need later.\nThis property relates $\\d$-refinements of partitions and $\\langle \\d \\rangle$-regularity of partitions,\nand relies\non Claim~\\ref{claim:refinement-union}. Here, as well as in the rest of this section, we will use\nthe definitions and notations introduced in Section \\ref{sec:define}.\nIn particular, recall that if a vertex partition $\\Z$ of vertex classes $(V_1,V_2,V_3)$ satisfies $\\Z \\prec \\{V_1,V_2,V_3\\}$,\nthen for every $1 \\le i \\le 3$ we denote $\\Z_i = \\{Z \\in \\Z \\,\\vert\\, Z \\sub V_i\\}$.\nMoreover, if a $2$-partition $(\\Z,\\E)$, satisfies $\\Z \\prec \\{V_1,V_2,V_3\\}$ we denote $\\E_i = \\{E \\in \\E \\,\\vert\\, E \\sub V_j \\times V_k\\}$ where $\\{i,j,k\\}=\\{1,2,3\\}$. We will first need the following easy claim regarding the union of $\\langle \\d\\rangle$-regular graphs.\n\n\n\\begin{claim}\\label{claim:star-union\n\tLet $G_1,\\ldots,G_\\ell$ be mutually edge-disjoint bipartite graphs on the same vertex classes $(Z,Z')$.\n\tIf every $G_i$ is $\\langle \\d \\rangle$-regular then $G=\\bigcup_{i=1}^\\ell G_i$ is also $\\langle \\d \\rangle$-regular.\n\\end{claim}\n\\begin{proof\n\n\tLet $S \\sub Z$, $S' \\sub Z'$ with $|S| \\ge \\d|Z|$, $|S'| \\ge \\d|Z'|$.\n\tThen\n\t$$d_G(S,S') = \\frac{e_G(S,S')}{|S||S'|}\n\t= \\sum_{i=1}^\\ell \\frac{e_{G_i}(S,S')}{|S||S'|}\n\t= \\sum_{i=1}^\\ell d_{G_i}(S,S')\n\t\\ge \\sum_{i=1}^\\ell \\frac12 d_{G_i}(Z,Z')\n\t= \\frac12 d_{G}(Z,Z') \\;,$$\n\twhere the second and last equalities follow from the mutual disjointness of the $G_i$, and the inequality follows from the $\\langle \\d \\rangle$-regularity of each $G_i$.\n\tThus, $G$ is $\\langle \\d \\rangle$-regular, as claimed.\n\\end{proof}\n\nWe use the following claim regarding approximate refinements.\n\n\\begin{claim}\\label{claim:refinement-union}\n\tIf $\\Q \\prec_\\d \\P$ then there exist $P \\in \\P$ and $Q$ that is a union of members of $\\Q$ such that $|P \\triangle Q| \\le 3\\d|P|$.\n\\end{claim}\n\\begin{proof}\n\tFor each $P\\in \\P$ let $\\Q(P) = \\{Q \\in \\Q \\colon Q \\sub_\\d P\\}$,\n\tand denote $P_\\Q = \\bigcup_{Q \\in \\Q(P)} Q$.\n\tWe have\n\t\\begin{align*}\n\t\\sum_{P \\in \\P} |P \\triangle P_\\Q|\n\t&= \\sum_{P \\in \\P} |P_\\Q \\sm P| + \\sum_{P \\in \\P} |P \\sm P_\\Q|\n\t= \\sum_{P \\in \\P} \\sum_{\\substack{Q \\in \\Q \\colon\\\\Q \\sub_\\d P}} |Q \\sm P|\n\t+ \\sum_{P \\in \\P} \\sum_{\\substack{Q \\in \\Q \\colon\\\\Q \\nsubseteq_\\d P}} |Q \\cap P| \\\\\n\t&\\le \\sum_{P \\in \\P} \\sum_{\\substack{Q \\in \\Q \\colon\\\\Q \\sub_\\d P}} \\d|Q|\n\t+ \\Big( \\sum_{\\substack{Q \\in \\Q\\colon\\\\Q \\notin_\\d \\P}} |Q|\n\t+ \\sum_{\\substack{Q \\in \\Q \\colon\\\\Q \\in_\\d \\P}} \\d|Q| \\Big)\n\t\\le 3\\d\\sum_{Q \\in \\Q} |Q|\n\t= 3\\d\\sum_{P \\in \\P} |P| \\;,\n\t\\end{align*}\n\twhere the last inequality uses the statement's assumption $\\Q \\prec_\\d \\P$ to bound the middle summand.\n\tBy averaging, there exists $P \\in \\P$ such that $|P \\triangle P_\\Q| \\le 3\\d|P|$, thus completing the proof.\n\\end{proof}\n\nThe property of $2$-partitions that we need is as follows.\n\\begin{claim}\\label{claim:uniform-refinement}\n\tLet $\\P=(\\Z,\\E)$ be a $2$-partition with $\\Z \\prec \\{V_1,V_2,V_3\\}$,\n\tand let $\\G$ be a partition of $V_1\\times V_2$ with $\\E_3 \\prec_\\d \\G$.\n\tIf $(\\Z,\\E)$ is $\\langle \\d \\rangle$-good\n\tthen $\\Z_1 \\cup \\Z_2$ is a $\\langle 3\\d \\rangle$-regular partition of some $G \\in \\G$.\n\\end{claim}\n\n\\begin{proof}\n\tPut $\\E=\\E_3$.\n\tBy Claim~\\ref{claim:refinement-union}, since $\\E \\prec_\\d \\G$ there exist $G \\in \\G$ (a bipartite graph on $(V_1,V_2)$) and $G_\\E$ that is a union of members of $\\E$ (and thus also a bipartite graph on $(V_1,V_2)$) such that $|G \\triangle G_\\E| \\le 3\\d|G|$.\n\n\tLetting $Z_1 \\in \\Z_1$, $Z_2 \\in \\Z_2$, to complete the proof it suffices to show that the induced bipartite graph $G_\\E[Z_1,Z_2]$ is $\\langle \\d \\rangle$-regular (recall Definition~\\ref{def:star-regular}).\n\tBy Definition~\\ref{def:2-partition}, $G_\\E[Z_1,Z_2]$ is a union of bipartite graphs from $\\E$ on $(Z_1,Z_2)$.\n\tSince every graph in $\\E$ is $\\langle \\d \\rangle$-regular by the statement's assumption that $(\\Z,\\E)$ is $\\langle \\d \\rangle$-good (recall Definition~\\ref{def:k-good}), we have that $G_\\E[Z_1,Z_2]$ is a union of $\\langle \\d \\rangle$-regular bipartite graphs on $(Z_1,Z_2)$.\n\tBy Claim~\\ref{claim:star-union}, $G_\\E[Z_1,Z_2]$ is $\\langle \\d \\rangle$-regular as well, thus completing the proof.\n\\end{proof}\n\nWe will later need the following easy (but slightly tedious to state) claim.\n\\begin{claim}\\label{claim:restriction}\n\tLet $H$ be a $3$-partite $3$-graph on vertex classes $(V_1,V_2,V_3)$, and let $H'$ be the induced $3$-partite $3$-graph on vertex classes $(V_1',V_2',V_3')$ with $V_i' \\sub V_i$ and $\\a \\cdot e(H)$ edges.\n\tIf $(\\Z,\\E)$ is a $\\langle \\d \\rangle$-regular partition of $H$ with\n\t$\\Z \\prec \\bigcup_{i=1}^3 \\{V_i,\\,V_i \\sm V_i'\\}$\n\n\tthen its restriction $(\\Z',\\E')$ to $V(H')$ is a $\\langle \\d\/\\a \\rangle$-regular partition of $H'$.\n\\end{claim}\n\\begin{proof}\n\tRecall Definition~\\ref{def:k-reg}.\n\tClearly, $(\\Z',\\E')$ is $\\langle \\d \\rangle$-good. We will show that $\\E'_1 \\cup \\Z'_1$ is a $\\langle \\d\/\\a \\rangle$-regular partition of $G^1_{H'}$.\n\tThe argument for $G^2_{H'}$, $G^3_{H'}$ will be analogous, hence the proof would follow.\n\tObserve that $G^1_{H'}$ is an induced subgraph of $G^1_{H}$, namely, $G^1_{H'} = G^i_{H}[V_2' \\times V_3',\\, V_1']$.\n\tBy assumption, $e(H') = \\a e(H)$, and thus $e(G^1_{H'}) = \\a e(G^1_{H})$.\n\tBy the statement's assumption on $\\Z$ and since $\\E_1 \\cup \\Z_1$ is a $\\langle \\d \\rangle$-regular partition of $G^1_{H}$, we deduce---by adding\/removing at most $\\d e(G^1_{H}) = (\\d\/\\a)e(G^1_{H'})$ edges of $G^1_{H'}$---that $\\E'_1 \\cup \\Z'_1$ is a $\\langle \\d\/\\a \\rangle$-regular partition of $G_{H'}^1$.\n\tAs explained above, this completes the proof.\n\\end{proof}\n\n\nFinally, we will need the following claim regarding approximate refinements.\n\\begin{claim}\\label{claim:refinement-size}\n\tIf $\\Q \\prec_{1\/2} \\P$ and $\\P$ is equitable then $|\\Q| \\ge \\frac14|\\P|$.\n\\end{claim}\n\\begin{proof}\n\tWe claim that the underlying set $U$ has a subset $U^*$ of size $|U^*|\\ge \\frac14|U|$ such that the partitions $\\Q^*=\\{Q \\cap U^* \\,\\vert\\, Q \\in \\Q \\} \\setminus \\{\\emptyset\\}$ and $\\P^*=\\{P \\cap U^* \\,\\vert\\, P \\in \\P \\} \\setminus \\{\\emptyset\\}$\n\n\tof $U^*$ satisfy $\\Q^* \\prec \\P^*$.\n\tIndeed, let $U^* = \\bigcup_{Q} Q \\cap P_Q$ where the union is over all $Q \\in \\Q$ satisfying $Q \\sub_{1\/2} P_Q$ for a (unique) $P_Q \\in \\P$.\n\n\tAs claimed, $|U^*| = \\sum_{Q \\in_{1\/2} \\P} |Q \\cap P_Q| \\ge \\sum_{Q \\in_{1\/2} \\P} \\frac12|Q| \\ge \\frac14|U|$, using $\\Q \\prec_{1\/2} \\P$ for the last inequality.\n\n\tNow, since $\\P$ is equitable, $|\\P^*| \\ge \\frac14|\\P|$.\n\n\tThus, $|\\Q| \\ge |\\Q^*| \\ge |\\P^*| \\ge \\frac14|\\P|$, as desired.\n\\end{proof}\n\n\n\\renewcommand{\\K}{k}\n\\renewcommand{\\w}{w}\n\n\\subsection{$3$-graph key argument}\\label{subsec:key}\n\n\nWe next introduce a few more definitions that are needed for the statement of Lemma \\ref{lemma:ind-k}.\nLet $e(i) = 2^{i+10}$. We define the following tower-type function $\\t\\colon\\N\\to\\N$;\n\\begin{equation}\\label{eq:t}\n\\t(i+1) = \\begin{cases}\n2^{\\t(i)\/e(i)}\t&\\text{if } i \\ge 1\\\\\n2^{250}\t&\\text{if } i = 0 \\;.\n\\end{cases}\n\\end{equation}\nIt is easy to prove, by induction on $i$, that $\\t(i) \\ge e(i)\\t(i-1)$ for $i \\ge 2$ (for the induction step,\n$t(i+1) \\ge 2^{\\t(i-1)} = t(i)^{e(i-1)}$, so $t(i+1)\/e(i+1) \\ge \\t(i)^{e(i-1)-i-11} \\ge \\t(i)$).\nThis means that $t$ is monotone increasing, and that $\\t$ is an integer power of $2$ (follows by induction as $t(i)\/e(i) \\ge 1$ is a positive power of $2$ and in particular an integer).\nWe record the following facts regarding $\\t$ for later use:\n\\begin{equation}\\label{eq:monotone}\n\\t(i) \\ge 4\\t(i-1) \\quad\\text{ and }\\quad\n\\text{ $\\t(i)$ is a power of $2$} \\;.\n\\end{equation}\nFor a function $f:\\N\\to\\N$ with $f(i) \\ge i$ we denote\n\\begin{equation}\\label{eq:f*}\nf^*(i) = \\t\\big(f(i)\\big)\/e(i) \\;.\n\\end{equation}\nNote that $f^*(i)$ is indeed a positive integer (by the monotonicity of $\\t$, $f^*(i) \\ge \\t(i)\/e(i)$ is a positive power of $2$).\nIn fact, $f^*(i) \\ge f(i)$ (as $f^*(i) \\ge 4^{f(i)}\/e(i)$ using~(\\ref{eq:monotone})).\nWe recursively define the function $\\w\\colon\\N\\to\\N$ as follows;\n\\begin{equation}\\label{eq:Ak}\n\\w(i+1) = \\begin{cases}\n\\w^*(i)\t\t&\\text{if } i \\ge 1\\\\\n1\t\t&\\text{if } i = 0 \\;.\n\\end{cases}\n\\end{equation}\nIt is evident that $\\w$ is a wowzer-type function; in fact, one can check that:\n\\begin{equation}\\label{eq:A_k}\n\\w(i) \\ge \\wow(i) \\;.\n\\end{equation}\n\n\n\n\n\n\n\n\n\\begin{lemma}[Key argument]\\label{lemma:ind-k}\n\n\n\tLet $s \\in \\N$, let $\\Vside^1,\\Vside^2,\\Vside^3$ be mutually disjoint sets of equal size and let $\\V^1 \\succ\\cdots\\succ \\V^m$ be a sequence of\n\t$m=\\w^*(s)+1$ successive equitable refinements of $\\{\\Vside^1,\\Vside^2,\\Vside^3\\}$\n\twith $|\\V^i_1|=|\\V^i_2|=|\\V^i_3|=\\t(i)$ for every\\footnote{Since we assume that each $\\V^i$ refines $\\{\\Vside^1,\\Vside^2,\\Vside^3\\}$ then $\\V^i_1$ is (by the notation mentioned before Claim \\ref{claim:star-union}) the restriction of $\\V^i$ to $\\Vside^1$.} $1 \\leq i \\leq m$.\n\tThen there is a $3$-partite $3$-graph $H$ on $(\\Vside^1,\\Vside^2,\\Vside^3)$ of density $d(H)=2^{-s}$ satisfying the following property:\\\\\n\tIf $(\\Z,\\E)$ is a $\\langle 2^{-70} \\rangle$-regular partition of $H$ and for some\n\t$1 \\le i \\le \\w(s)$ $(< m)$ we have $\\Z_3 \\prec_{2^{-9}} \\V^i_3$ and $\\Z_2 \\prec_{2^{-9}} \\V^i_2$ then we also have $\\Z_1 \\prec_{2^{-9}} \\V^{i+1}_1$.\n\\end{lemma}\n\n\\begin{proof}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\tPut $s':=\\w^*(s)$, so that $m = s'+1$.\n\tApply Lemma~\\ref{theo:core} with\n\t$$\\Lside=\\Vside^1,\\quad \\Rside=\\Vside^2 \\quad\\text{ and }\\quad \\V^2_1 \\succ \\cdots \\succ \\V^{s'+1}_1 ,\\quad \\V^1_2 \\succ \\cdots \\succ \\V^{s'}_2 \\;,$$\n\tand let\n\t\\begin{equation}\\label{eq:main-k-colors}\n\t\\G^1 \\succ \\cdots \\succ \\G^{s'} \\quad\\text{ with }\\quad |G^\\ell|=2^\\ell \\text{ for every } 1 \\le \\ell \\le s'\n\t\\end{equation}\n\tbe the resulting sequence of $s'$ successively refined equipartitions of $\\Vside^1 \\times \\Vside^2$.\n\n\n\n\n\t\n\t\\begin{prop}\\label{prop:main-k-hypo}\n\t\tLet $1 \\le \\ell \\le s'$ and $G \\in \\G^\\ell$.\n\t\tFor every $\\langle 2^{-28} \\rangle$-regular partition $\\Z_1 \\cup \\Z_2$ of $G$ (where $\\Z_1$ and $\\Z_2$ are partitions of $\\Vside^1$ and $\\Vside^2$, respectively) and every $1 \\le i \\le \\ell$,\n\t\tif $\\Z_2 \\prec_{2^{-9}} \\V^i_2$\n\t\tthen $\\Z_1 \\prec_{2^{-9}} \\V^{i+1}_1$.\n\t\\end{prop}\n\t\\begin{proof}\n\t\tFirst we need to verify that we may apply Lemma~\\ref{theo:core} as above.\n\t\tAssumptions~\\ref{item:core-minR},~\\ref{item:core-expR} in Lemma~\\ref{theo:core} hold by~(\\ref{eq:monotone}) and the fact that $|\\V^j_2|=\\t(j)$.\n\t\tAssumption~\\ref{item:core-expL} is satisfied since for every $1 \\le j \\le s$ we have\n\t\t$$|\\V^{j+1}_1| = \\t(j+1) = 2^{\\t(j)\/e(j)} = 2^{|\\V^{j}_2|\/e(j)} \\;,$$\n\t\twhere the second equality uses the definition of the function $\\t$ in~(\\ref{eq:t}).\n\t\tWe can thus use Lemma~\\ref{theo:core} to infer that the fact that $\\Z_2 \\prec_{2^{-9}} \\V^i_2$ implies that $\\Z_1 \\prec_x \\V^{i+1}_1$ with $x=\\max\\{2^{5}\\sqrt{2^{-28}},\\, 32\/\\sqrt[6]{\\t(1)} \\} = 2^{-9}$, using~(\\ref{eq:t}).\n\t\\end{proof}\n\t\n\n\n\n\n\n\n\n\n\n\n\t\n\tFor each $1 \\le j \\le s$ let\n\t\\begin{equation}\\label{eq:main-k-dfns}\n\t\\G^{(j)} = \\G^{\\w^*(j)}\n\t\\quad\\text{ and }\\quad\n\t\\V^{(j)} = \\V^{\\w(j)}_3 \\;.\n\t\\end{equation}\t\n\tAll these choices are well defined since $\\w^*(j)$ satisfies $1 \\le \\w^*(1) \\le \\w^*(j) \\le \\w^*(s) = s'$, and since $\\w(j)$ satisfies $1 \\le \\w(1) \\le \\w(j) \\le \\w(s) \\le m$. Observe that we have thus chosen two subsequences of $\\G^1,\\cdots,\\G^{s'}$ and $\\V^1_3,\\ldots,\\V^m_3$, each of length $s$.\n\tRecalling that each $\\G^{(j)}$ is a partition of $\\Vside^1 \\times \\Vside^2$, we now apply Lemma~\\ref{theo:core} again with\n\t$$\n\t\\Lside=\\Vside^1 \\times \\Vside^2,\\quad \\Rside=\\Vside^3 \\quad\\text{ and }\\quad \\G^{(1)} \\succ \\cdots \\succ \\G^{(s)}, \\quad \\V^{(1)} \\succ \\cdots \\succ \\V^{(s)} \\;.\n\t$$\n\tThe output of this application of Lemma~\\ref{theo:core} consists of a sequence of $s$ (successively refined)\n\tequipartitions of $(\\Vside^1 \\times \\Vside^2)\\times\\Vside^3$.\n\tWe can think of the $s$-th partition of this sequence as a collection of $2^s$ bipartite graphs on vertex sets\n\t$(\\Vside^1\\times\\Vside^{2},\\,\\Vside^3)$. For the rest of the proof let $G'$ be be any of these graphs.\n\tBy Remark \\ref{remequi} we have\n\t\\begin{equation}\\label{eq:ind-colors2}\n\td(G')=2^{-s} \\;.\n\t\\end{equation}\n\t\\begin{prop}\\label{prop:ind-prop2}\n\t\tFor every $\\langle 2^{-70} \\rangle$-regular partition $\\E \\cup \\V$ of $G'$ (where $\\E$ and $\\V$ are partitions of $\\Vside^1\\times\\Vside^{2}$ and $\\Vside^3$ respectively)\n\t\tand every $1 \\le j' \\le s$,\n\t\tif $\\V \\prec_{2^{-9}} \\V^{(j')}$ then $\\E \\prec_{2^{-30}} \\G^{(j')}$.\n\t\\end{prop}\n\t\\begin{proof}\n\t\tFirst we need to verify that we may apply Lemma~\\ref{theo:core} as above.\n\t\tNote that $|\\G^{(j)}|=2^{\\w^*(j)}$ by~(\\ref{eq:main-k-colors}) and (\\ref{eq:main-k-dfns}),\n\t\tand that $|\\V^{(j)}|=\\t(\\w(j))$ by (\\ref{eq:main-k-dfns}) and the statement's assumption that $|\\V^i_3|=\\t(i)$.\n\t\tTherefore,\n\t\t\\begin{equation}\\label{eq:main-k-orders}\n\t\t|\\G^{(j)}| = 2^{\\w^*(j)} = 2^{\\t(\\w(j))\/e(j)} = 2^{|\\V^{(j)}|\/e(j)} \\;,\n\t\n\t\t\\end{equation}\n\t\twhere the second equality relies on~(\\ref{eq:f*}).\n\t\tMoreover, note that $\\t(\\w(1)) = \\t(1) = 2^{300}$.\n\t\tNow, Assumptions~\\ref{item:core-minR} and~\\ref{item:core-expR} in Lemma~\\ref{theo:core} follow from the fact that $|\\V^{(j)}|=\\t(\\w(j))$, from~(\\ref{eq:monotone})\n\t\tand the fact that $|\\V^{(1)}| = \\t(\\w(1)) \\ge 2^{200}$ by~(\\ref{eq:Ak}). \t\n\t\n\t\tAssumption~\\ref{item:core-expL} follows from~(\\ref{eq:main-k-orders}).\n\t\tWe can thus use Lemma~\\ref{theo:core} to infer that the fact that $\\V \\prec_{2^{-9}} \\V^{(j')}$ implies that\n\t\t$\\E \\prec_x \\G^{(j')}$ with $x=\\max\\{2^{5}\\sqrt{2^{-70}},\\, 32\/\\sqrt[6]{\\t(\\w(1))} \\} = 2^{-30}$.\n\t\n\t\n\t\n\t\\end{proof}\n\t\n\tLet $H$ be the $3$-partite $3$-graph on vertex classes $(\\Vside^1,\\Vside^2,\\Vside^3)$ with edge set\n\t$$\n\tE(H) = \\big\\{ (v_1,v_2,v_3) \\,:\\, ((v_1,v_{2}),v_3) \\in E(G') \\big\\} \\;,\n\t$$\t\n\tand note that we have (recall Definition \\ref{def:aux})\n\t\\begin{equation}\\label{eqH}\n\tG'=G_{H}^3\\;.\n\t\\end{equation}\n\tWe now prove that $H$ satisfies the properties in the statement of the lemma.\n\n\n\n\n\tFirst, note that by~(\\ref{eq:ind-colors2}) and (\\ref{eqH}) we have $d(H)=2^{-s}$, as needed.\n\tAssume now that $i$ is such that\n\t\\begin{equation}\\label{eq:ind-i-assumption}\n\t1 \\le i \\le \\w(s)\n\t\\end{equation}\n\tand:\n\t\\begin{enumerate}\n\t\t\\item\\label{item:ind-reg}\n\t\t$(\\Z,\\E)$ is a $\\langle 2^{-70} \\rangle$-regular partition of $H$, and\n\t\t\\item\\label{item:ind-refine} $\\Z_3 \\prec_{2^{-9}} \\V^i_3$ and $\\Z_2 \\prec_{2^{-9}} \\V^i_2$.\n\t\\end{enumerate}\t\n\tWe need to show that\n\t\\begin{equation}\\label{eq:ind-goal}\n\t\\Z_1 \\prec_{2^{-9}} \\V^{i+1}_1 \\;.\n\t\\end{equation}\n\n\tSince Item~\\ref{item:ind-reg} guarantees that $(\\Z,\\E)$ is a $\\langle 2^{-70} \\rangle$-regular partition of $H$,\n\n\twe get from Definition~\\ref{def:k-reg} and (\\ref{eqH}) that\n\t\\begin{equation}\\label{eq:ind-reg}\n\t\\text{$\\E_3 \\cup \\Z_3$ is a $\\langle 2^{-70} \\rangle$-regular partition of } G'.\n\t\\end{equation}\n\tLet\n\t\\begin{equation}\\label{eq:ind-j'}\n\t1 \\le j' \\le s\n\t\\end{equation}\n\tbe the unique integer satisfying (the equality here is just (\\ref{eq:Ak}))\n\t\\begin{equation}\\label{eq:ind-sandwich}\n\t\\w(j') \\le i < \\w(j'+1) = \\w^*(j')\\;.\n\t\\end{equation}\n\tNote that (\\ref{eq:ind-j'}) holds due to~(\\ref{eq:ind-i-assumption}).\n\tRecalling~(\\ref{eq:main-k-dfns}),\n\tthe lower bound in~(\\ref{eq:ind-sandwich}) implies that $\\V^i_3 \\prec \\V^{\\w(j')} = \\V^{(j')}$.\n\tTherefore, the assumption $\\Z_3 \\prec_{2^{-9}} \\V^i_3$ in~\\ref{item:ind-refine} implies that\n\t\\begin{equation}\\label{eq:ind-Zk}\n\t\\Z_3 \\prec_{2^{-9}} \\V^{(j')} \\;.\n\t\\end{equation}\n\tApply Proposition~\\ref{prop:ind-prop2} on $G'$, using~(\\ref{eq:ind-reg}),~(\\ref{eq:ind-j'}) and~(\\ref{eq:ind-Zk}), to deduce that\n\t\\begin{equation}\\label{eq:ind-E}\n\t\\E_3 \\prec_{2^{-30}} \\G^{(j')} = \\G^{\\w^*(j')} \\;,\n\t\\end{equation}\n\twhere for the equality again recall~(\\ref{eq:main-k-dfns}).\n\tSince $(\\Z,\\E)$ is a $\\langle 2^{-70} \\rangle$-regular partition of $H$ (by Item~\\ref{item:ind-reg} above)\n\tit is in particular $\\langle 2^{-70} \\rangle$-good. By~(\\ref{eq:ind-E}) we may thus apply Claim~\\ref{claim:uniform-refinement}\n\tto conclude that\n\t\\begin{equation}\\label{eq:ind-reg2}\n\t\\Z_1 \\cup \\Z_2 \\text{ is a }\n\t\\langle 2^{-28} \\rangle\n\t\\text{-regular partition of some $G\\in\\G^{\\w^*(j')}$.}\n\t\\end{equation}\n\tBy~(\\ref{eq:ind-reg2}) we may apply Proposition~\\ref{prop:main-k-hypo} with $G$, $\\Z_1\\cup\\Z_2$, $\\ell=\\w^*(j')$ and $i$, observing (crucially)\n\tthat $i \\leq \\ell$ by (\\ref{eq:ind-sandwich}). We thus conclude that the fact $\\Z_2 \\prec_{2^{-9}} \\V^i_2$ (stated in~\\ref{item:ind-refine})\n\timplies that $\\Z_1 \\prec_{2^{-9}} \\V^{i+1}_1$, thus proving~(\\ref{eq:ind-goal}) and completing the proof.\n\\end{proof}\n\n\n\n\n\\subsection{Putting everything together}\\label{subsec:pasting}\n\n\nWe can now prove our main theorem, Theorem~\\ref{theo:main}, which we repeat here for convenience.\n\\addtocounter{theo}{-2}\n\\begin{theo}[Main theorem\n\tLet $s \\in \\N$.\n\tThere exists a $3$-partite $3$-graph $H$ on vertex classes of equal size and of density at least $2^{-s}$,\n\tand a partition $\\V_0$ of $V(H)$ with $|\\V_0| \\le 2^{300}$, such that if\n\t$(\\Z,\\E)$ is a $\\langle 2^{-73} \\rangle$-regular partition of $H$ with\n\t$\\Z \\prec \\V_0$ then $|\\Z| \\ge \\wow(s)$.\n\n\n\n\n\n\n\n\n\n\n\n\t\n\n\n\n\n\n\n\n\n\n\n\n\\end{theo}\n\\addtocounter{theo}{+1}\n\n\n\n\\begin{proof}\n\n\n\n\n\n\n\n\tLet the $3$-graph $B$ be the tight $6$-cycle; that is, $B$ is the $3$-graph on vertex classes $\\{0,1,\\ldots,5\\}$ with edge set $E(B)=\\{\\{0,1,2\\},\\{1,2,3\\},\\{2,3,4\\},\\{3,4,5\\},\\{4,5,0\\},\\{5,0,1\\}\\}$.\n\tNote that $B$ is $3$-partite with vertex classes $(\\{0,3\\},\\{1,4\\},\\{2,5\\}\\}$.\n\n\tPut $m=\\w^*(s-1)+1$ and let $n \\ge \\t(m)$.\n\tLet $\\Vside^0,\\ldots,\\Vside^{5}$ be $6$ mutually disjoint sets of size $n$ each.\n\tLet $\\V^1 \\succ\\cdots\\succ \\V^m$ be an arbitrary sequence of $m$ successive equitable refinements of $\\{\\Vside^0,\\ldots,\\Vside^{5}\\}$ with $|\\V^i_h|=\\t(i)$ for every $1 \\le i \\le m$ and $0 \\le h \\le 5$, which exists as $n$ is large enough.\n\n\tExtending the notation $\\Z_i$ (above Definition~\\ref{def:k-reg}), for every $0 \\le x \\le 5$\n\n\twe henceforth denote the restriction of the vertex partition $\\Z$ to $\\Vside^x$ by $\\Z_x = \\{Z \\in \\Z \\,\\vert\\, Z \\sub \\Vside^x\\}$.\t\n\tFor each edge $e=\\{x,x+1,x+2\\} \\in E(B)$ (here and henceforth when specifying an edge, the integers are implicitly taken modulo $6$)\n\tapply Lemma~\\ref{lemma:ind-k} with\n\t$$s-1,\\,\n\t\\Vside^{x},\\Vside^{x+1},\\Vside^{x+2}\n\t\\text{ and }\n\t(\\V^{1}_x \\cup \\V^1_{x+1} \\cup \\V^1_{x+2})\n\t\\succ\\cdots\\succ (\\V^{m}_{x}\\cup\\V^{m}_{x+1}\\cup\\V^{m}_{x+2}) \\;.$$\n\n\tLet $H_e$ denote the resulting $3$-partite $3$-graph on $(\\Vside^{x},\\Vside^{x+1},\\Vside^{x+2})$.\n\tNote that $d(H_e) = 2^{-(s-1)}$.\n\tMoreover, let\n\t$$c = 2^{-9} \\quad\\text{ and }\\quad K=\\w(s-1)+1 \\;.$$\n\tThen $H_e$ has the property that for every $\\langle 2^{-70} \\rangle$-regular partition $(\\Z',\\E')$ of $H_e$ and every $1 \\le i < K$,\n\n\n\n\n\n\n\n\t\\begin{equation}\\label{eq:paste-property}\n\t\\text{if $\\Z'_{x+2} \\prec_{c} \\V^i_{x+2}$ and $\\Z'_{x+1} \\prec_{c} \\V^i_{x+1}$ then $\\Z'_x \\prec_{c} \\V^{i+1}_x$.}\n\t\\end{equation}\t\n\tWe construct our $3$-graph on the vertex set $\\Vside:=\\Vside^0 \\cup\\cdots\\cup \\Vside^5$ as\n\t$E(H) = \\bigcup_{e} E(H_e)$; that is, $H$ is the edge-disjoint union of all six $3$-partite $3$-graphs $H_e$ constructed above.\n\tNote that $H$ is a $3$-partite $3$-graph (on vertex classes $(\\Vside^0 \\cup \\Vside^3,\\, \\Vside^1 \\cup \\Vside^4,\\, \\Vside^2 \\cup \\Vside^5))$ of density $\\frac68 2^{-(s-1)} \\ge 2^{-s}$, as needed.\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\tWe will later use the following fact.\n\t\\begin{prop}\\label{prop:restriction}\n\t\tLet $(\\Z,\\E)$ be an $\\langle 2^{-73}\\rangle$-regular partition of $H$ and let\n\t\t$e \\in E(B)$.\n\t\n\t\n\t\n\t\tIf $\\Z \\prec \\{\\Vside^0,\\ldots,\\Vside^{5}\\}$\n\t\n\t\tthen the restriction $(\\Z',\\E')$ of $(\\Z,\\E)$ to $V(H_e)$ is a $\\langle 2^{-70} \\rangle$-regular partition of $H_e$.\n\t\\end{prop}\n\t\\begin{proof}\n\t\tImmediate from Claim~\\ref{claim:restriction} using the fact that $e(H_e) = \\frac16 e(H)$.\n\t\n\t\\end{proof}\t\n\n\t\n\tNow, let $(\\Z,\\E)$ be a $\\langle 2^{-73} \\rangle$-regular partition of $H$\n\twith $\\Z \\prec \\V^1$.\n\tOur goal will be to show that\n\t\\begin{equation}\\label{eq:paste-goal}\n\t\\Z \\prec_{c} \\V^{K} \\;.\n\t\\end{equation}\n\tProving~(\\ref{eq:paste-goal}) would complete the proof, by setting $\\V_0$ in the statement to be $\\V^1$ here (notice $|\\V^1|=3\\t(1) \\le 2^{300}$ by~(\\ref{eq:t}));\n\tindeed, Claim~\\ref{claim:refinement-size} would imply that\n\t$$|\\Z| \\ge \\frac14|\\V^{K}| = \\frac14 \\cdot 6 \\cdot \\t(K)\n\t\\ge \\t(K)\n\t\\ge \\t(\\w(s-1))\n\t\\ge \\w(s)\n\t\\ge \\wow(s) \\;,$$\n\n\n\twhere the last inequality uses~$(\\ref{eq:A_k})$.\n\n\n\n\tAssume towards contradiction that $\\Z \\nprec_{c} \\V^{K}$. By averaging,\n\t\\begin{equation}\\label{eq:assumption}\n\t\\Z_h \\nprec_c \\V^{K}_h \\text{ for some } 0 \\le h \\le 5.\n\t\\end{equation}\n\tFor each $0 \\le h \\le 5$ let $1 \\le \\b(h) \\le K$ be the largest integer satisfying $\\Z_h \\prec_c \\V^{\\b(h)}_h$,\n\twhich is well defined since $\\Z_h \\prec_c \\V^1_h$,\n\tsince in fact $\\Z \\prec \\V^1$.\n\tPut $\\b^* = \\min_{0 \\le h \\le 5} \\b(h)$, and note that by~(\\ref{eq:assumption}),\n\t\\begin{equation}\\label{eq:paste-star}\n\t\\b^* < K \\;.\n\t\\end{equation}\n\tLet $0 \\le x \\le 5$ minimize $\\b$, that is, $\\b(x)=\\b^*$.\n\tTherefore:\n\t\\begin{equation}\\label{eqcontra}\n\t\\Z_{x+2} \\prec_c \\V^{\\b^*}_{x+2} \\mbox{~,~} \\Z_{x+1} \\prec_c \\V^{\\b^*}_{x+1} \\mbox{ and }\n\t\\Z_{x} \\nprec_c \\V^{\\b^*+1}_{x}.\t\t\n\t\\end{equation}\n\tLet $e=\\{x,x+1,x+2\\} \\in E(B)$.\n\tLet $(\\Z',\\E')$ be the restriction of $(\\Z,\\E)$ to $V(H_e)=\\Vside^{x} \\cup \\Vside^{x+1} \\cup \\Vside^{x+2}$, which is a $\\langle 2^{-70} \\rangle$-regular partition of $H_e$ by Proposition~\\ref{prop:restriction}. Since $\\Z'_x=\\Z_x$, $\\Z'_{x+1}=\\Z_{x+1}$, $\\Z'_{x+2}=\\Z_{x+2}$ we get\n\tfrom (\\ref{eqcontra}) a contradiction to~(\\ref{eq:paste-property}) with $i=\\beta^*$.\n\tWe have thus proved~(\\ref{eq:paste-goal}) and so the proof is complete.\n\\end{proof}\n\n\n\n\t\n\t\n\n\n\n\n\\section{Wowzer-type Lower Bounds for $3$-Graph Regularity Lemmas}\\label{sec:FR}\n\n\\renewcommand{\\K}{\\mathcal{K}}\n\nThe purpose of this section is to apply Theorem \\ref{theo:main} in order to prove Corollary~\\ref{coro:FR-LB},\nthus giving wowzer-type (i.e., $\\Ack_3$-type) lower bounds for the $3$-graph regularity lemmas of Frankl and R\\\"odl~\\cite{FrankRo02} and of Gowers~\\cite{Gowers06}.\nWe will start by giving the necessary definitions for Frankl and R\\\"odl's lemma\nand state our corresponding lower bound.\nNext we will state the necessary definitions for Gower's lemma and state our corresponding lower bound.\nThe formal proofs would then follow.\n\n\n\n\n\\subsection{Frankl and R\\\"odl's $3$-graph regularity}\n\n\n\n\n\n\n\\begin{definition}[$(\\ell,t,\\e_2)$-equipartition,~\\cite{FrankRo02}]\\label{def:ve-partition}\n\n\tAn \\emph{$(\\ell,t,\\e_2)$-equipartition} on a set $V$ is a $2$-partition $(\\Z,\\E)$ on $V$ where $\\Z$ is an equipartition of order $|\\Z|=t$ and every graph in $\\E$ is $\\e_2$-regular\\footnote{Here, and in several places in this section, we of course refer to the ``traditional'' notion of Szemer\\'edi's $\\e$-regularity, as defined at the beginning of Section \\ref{sec:define}. } of density $\\ell^{-1} \\pm \\e_2$.\n\\end{definition}\n\n\\begin{remark}\n\tIf $\\e_2 \\le \\frac12\\ell^{-1}$ then $(Z,\\E)$ has at most $2\\ell$ bipartite graphs between every pair of clusters of $\\Z$.\n\\end{remark}\n\n\n\nA \\emph{triad} of a $2$-partition $(\\Z,\\E)$ is any tripartite graph whose three vertex classes are in $\\Z$ and three edge sets are in $\\E$. We often identify a triad with a triple of its edge sets $(E_1,E_2,E_3)$.\nThe \\emph{density} of a triad $P$ in a $3$-graph $H$ is $d_H(P)=|E(H) \\cap T(P)|\/|T(P)|$ (and $0$ if $|T(P)|=0$).\nA \\emph{subtriad} of $P$ is any subgraph of $P$ on the same vertex classes.\n\n\n\n\n\n\n\n\\begin{definition}[$3$-graph $\\e$-regularity~\\cite{FrankRo02}]\\label{def:FR-reg}\n\n\tLet $H$ be a $3$-graph.\n\tA triad $P$ is \\emph{$\\e$-regular} in $H$ if every subtriad $P'$ with $|T(P')| \\ge \\e|T(P)|$ satisfies $|d_H(P')-d_H(P)| \\le \\e$.\\\\\n\tAn $(\\ell,t,\\e_2)$-equipartition $\\P$ on $V(H)$ is an \\emph{$\\e$-regular} partition of $H$ if\n\n\t$\\sum_P |T(P)| \\le \\e|V|^3$ where the sum is over all triads of $\\P$ that are not $\\e$-regular in $H$.\t\n\\end{definition}\n\nThe $3$-graph regularity of Frankl and R\\\"odl~\\cite{FrankRo02} states, very roughly, that for every $\\e>0$ and every function $\\e_2\\colon\\N\\to(0,1]$, every $3$-graph has an $\\e$-regular $(\\ell,t,\\e_2(\\ell))$-equipartition where $t,\\ell$ are bounded by a wowzer-type function.\nIn fact, the statement in~\\cite{FrankRo02} uses a considerably stronger notion of regularity of a partition than in Definition~\\ref{def:FR-reg} that involves an additional function $r(t,\\ell)$ which we shall not discuss here (as discussed in~\\cite{FrankRo02}, this stronger notion was crucial for allowing them to prove the $3$-graph removal lemma).\nOur lower bound below applies even to the weaker notion stated above, which corresponds to taking $r(t,\\ell)\\equiv 1$.\n\n\n\nUsing Theorem~\\ref{theo:main} we can deduce a wowzer-type \\emph{lower} bound for Frankl and R\\\"odl's $3$-graph regularity lemma.\nThe proof of this lower bound appears in Subsection~\\ref{subsec:FR-LB-proof}.\n\n\n\n\\begin{theo}[Lower bound for Frankl and R\\\"odl's regularity lemma]\\label{theo:FR-LB}\n\tPut $c = 2^{-400}$.\n\tFor every $s \\in \\N$ there exists a $3$-partite $3$-graph $H$ of density $p=2^{-s}$, and a partition $\\V_0$ of $V(H)$ with $|\\V_0| \\le 2^{300}$,\n\tsuch that\n\tif $(\\Z,\\E)$ is an $\\e$-regular $(\\ell,t,\\e_2(\\ell))$-equipartition of $H$,\n\twith $\\e \\le c p$, $\\e_2(\\ell) \\le c \\ell^{-3}$ and $\\Z \\prec \\V_0$, then $|\\Z| \\ge \\wow(s)$.\n\\end{theo}\n\n\\begin{remark}\n\tOne can easily remove the assumption $\\Z \\prec \\V_0$ by taking the common refinement of $\\Z$ with $\\V_0$ (and adjusting $\\E$ appropriately).\nSince $|\\V_0|=O(1)$ this has only a minor effect on the parameters $\\e,\\ell,t,\\e_2(\\ell)$ of the partition and thus\none gets essentially the same lower bound. We omit the details of this routine transformation.\n\\end{remark}\n\n\n\\subsection{Gowers' $3$-graph regularity}\n\nHere we consider the $3$-graph regularity Lemma due to Gowers~\\cite{Gowers06}.\n\n\n\\begin{definition}[$\\a$-quasirandomness, see Definition~6.3 in \\cite{Gowers06}]\\label{def:quasirandom}\n\tLet $H$ be a $3$-graph, and let $P=(E_0,E_1,E_2)$ be a triad with $d(E_0)=d(E_1)=d(E_2)=:d$ on vertex classes $(X,Y,Z)$ with $|X|=|Y|=|Z|=:n$. We say that $P$ is \\emph{$\\a$-quasirandom} in $H$ if\n\t$$\\sum_{x_0,x_1 \\in X}\\sum_{y_0,y_1 \\in Y}\\sum_{z_0,z_1 \\in Z} \\prod_{i,j,k\\in\\{0,1\\}} f(x_i,y_j,z_k) \\le \\a d^{12}n^6 \\;,$$\n\n\twhere\n\t$$f(x,y,z) =\n\t\\begin{cases}\n\t1-d_H(P)\t\t\t\t&\\text{if } (x,y,z) \\in T(P), (x,y,z) \\in E(H)\\\\\n\t-d_H(P) \t\t\t\t&\\text{if } (x,y,z) \\in T(P), (x,y,z) \\notin E(H)\\\\\n\t0 \t\t\t\t\t\t&\\text{if } (x,y,z) \\notin T(P) \\;.\n\t\\end{cases}$$\n\tAn $(\\ell,t,\\e_2)$-equipartition $\\P$ on $V(H)$ is an \\emph{$\\a$-quasirandom} partition of $H$ if\n\n\t$\\sum_P |T(P)| \\le \\a|V|^3$ where the sum is over all triads of $\\P$ that are not $\\a$-quasirandom in $H$.\t\n\\end{definition}\n\nThe $3$-graph regularity lemma of Gowers~\\cite{Gowers06} (see also~\\cite{NaglePoRoSc09}) can be equivalently phrased as stating that, very roughly, for every $\\a>0$ and every function $\\e_2\\colon\\N\\to(0,1]$, every $3$-graph has an $\\a$-quasirandom $(\\ell,t,\\e_2(\\ell))$-equipartition where $t,\\ell$ are bounded by a wowzer-type function.\n\nOne way to prove a wowzer-type lower bound for Gowers' $3$-graph regularity lemma is along similar lines to the proof of Theorem~\\ref{theo:FR-LB}.\nHowever, there is shorter proof using the fact that Gowers' notion of quasirandomness implies Frankl and R\\\"odl's notion of regularity.\nIn all that follows we make the rather trivial assumption that, in the notation above, $\\a,1\/\\ell \\le 1\/2$.\n\n\\begin{prop}[\\cite{NagleRoSc17}]\\label{prop:Schacht}\n\tThere is $C \\ge 1$ such that the following holds;\n\tif a triad $P=(E_0,E_1,E_2)$ is $\\e^C$-quasirandom and for every $0 \\le i \\le 2$ the bipartite graph $E_i$ is $d(E_i)^C$-regular then $P$ is $\\e$-regular.\t\n\n\n\n\n\n\\end{prop}\n\n\nOur lower bound for Gowers' $3$-graph regularity lemma is as follows.\n\n\n\\begin{theo}[Lower bound for Gowers' regularity lemma]\\label{theo:Gowers-LB}\n\n\n\tFor every $s \\in \\N$ there exists a $3$-partite $3$-graph $H$ of density $p=2^{-s}$, and a partition $\\V_0$ of $V(H)$ with $|\\V_0| \\le 2^{300}$,\n\tsuch that\n\tif $(\\Z,\\E)$ is an $\\a$-quasirandom $(\\ell,t,\\e_2(\\ell))$-equipartition of $H$,\n\n\twith $\\a \\le \\poly(p)$, $\\e_2(\\ell) \\le \\poly(1\/\\ell)$ and $\\Z \\prec \\V_0$, then $|\\Z| \\ge \\wow(s)$.\n\\end{theo}\n\n\n\n\n\n\\begin{proof\n\n\tGiven $s$, let $H$ and $\\V_0$ be as in Theorem~\\ref{theo:FR-LB}.\n\tLet $\\P=(\\Z,\\E)$ be an $\\a$-quasirandom $(\\ell,t,\\e_2(\\ell))$-equipartition of $H$ with $\\Z \\prec \\V_0$, $\\a \\le (cp)^C$ and $\\e_2(\\ell) \\le \\min\\{c\\ell^{-3},\\,(2\\ell)^{-C}\\}$, where $c$ and $C$ are as in Theorem~\\ref{theo:FR-LB} and Proposition~\\ref{prop:Schacht} respectively.\n\tWe will show that $\\P$ is a $cp$-regular partition of $H$,\n\n\twhich would complete the proof using Theorem~\\ref{theo:FR-LB} and the fact that $\\e_2 \\le c\\ell^{-3}$.\n\tLet $P=(E_0,E_1,E_2)$ be a triad of $\\P$ that is $\\a$-quasirandom in $H$.\n\tNote that, by our choice of $\\e_2(\\ell)$, for every $0 \\le i \\le 2$ we have $d(E_i) \\ge 1\/\\ell - \\e_2(\\ell) \\ge 1\/2\\ell$; thus, since $\\e_2(\\ell) \\le (1\/2\\ell)^{C} \\le d(E_i)^C$, we have that $E_i$ is $d(E_i)^C$-regular.\n\tApplying Proposition~\\ref{prop:Schacht} on $P$ we deduce that $P$ is $\\e$-regular with $\\e=\\a^{1\/C} \\le cp$.\n\tSince $\\P$ is an $\\a$-quasirandom partition of $H$ we have, by Definition~\\ref{def:quasirandom} and since $\\a \\le \\e$, that $\\P$ is an $\\e$-regular partition of $H$, as needed.\n\n\n\\end{proof}\n\n\n\\subsection{Proof of Theorem~\\ref{theo:FR-LB}}\\label{subsec:FR-LB-proof}\n\nThe proof of Theorem~\\ref{theo:FR-LB} will follow quite easily from Theorem~\\ref{theo:main} together with Claim~\\ref{claim:reduction} below.\nClaim~\\ref{claim:reduction} basically shows that a $\\langle \\d \\rangle$-regularity ``analogue'' of Frankl and R\\\"odl's notion of regularity implies graph $\\langle \\d \\rangle$-regularity.\nHere it will be convenient to say that a graph partition is \\emph{perfectly} $\\langle \\d \\rangle$-regular if all pairs of distinct clusters are $\\langle \\d \\rangle$-regular without modifying any of the graph's edges.\nFurthermore, we will henceforth abbreviate $t(P)=|T(P)|$ for a triad $P$.\nWe will only sketch the proof of Claim~\\ref{claim:reduction}, deferring the full details to the Appendix~\\ref{sec:FR-appendix}.\n\n\\begin{claim}\\label{claim:reduction}\n\tLet $H$ be a $3$-partite $3$-graph on vertex classes $(\\Aside,\\Bside,\\Cside)$,\n\tand let $(\\Z,\\E)$ be an $(\\ell,t,\\e_2)$-equipartition of $H$ with $\\Z \\prec \\{\\Aside,\\Bside,\\Cside\\}$ such that for every triad $P$ of $\\P$ and every subtriad $P'$ of $P$ with $t(P') \\ge \\d \\cdot t(P)$ we have $d_H(P') \\ge \\frac23 d_H(P)$.\n\n\tIf $\\e_2(\\ell) \\le (\\d^2\/88)\\ell^{-3}$\n\tthen $\\E_3 \\cup \\Z_3$ is a perfectly $\\langle 2\\sqrt{\\d} \\rangle$-regular partition of $G_H^3$.\n\\end{claim}\n\n\\begin{proof}[Proof (sketch):]\nWe remind the reader that the vertex classes of $G_H^3$ are $(\\Aside \\times \\Bside,\\, \\Cside)$ (recall Definition~\\ref{def:aux}), and that $\\E_3$ and $\\Z_3$ are the partition of $\\Aside\\times\\Bside$ induced by $\\E$ and the partition of $\\Cside$ induced by $\\Z$, respectively. Suppose $(\\Z,\\E)$ is as in the statement of the claim, and define $\\E'$ as follows:\nfor every $A \\in \\Z_1$ and $C \\in \\Z_3$, replace all the bipartite graphs between $A$ and $C$ with the complete bipartite graph $A \\times C$.\nDo the same for every $B \\in \\Z_2$ and $C \\in \\Z_3$ (we do {\\em not} change the partitions between $\\Aside$ and $\\Bside$). The simple (yet somewhat tedious to prove)\nobservation is that if all triads of $(\\Z,\\E)$ are regular then all triads of $(\\Z,\\E')$ are essentially as regular.\nOnce this observation is proved, the proof of the claim reduces to checking definitions. We thus defer the proof to Appendix~\\ref{sec:FR-appendix}.\n\\end{proof}\n\n\nUsing Claim~\\ref{claim:reduction}, we now prove our wowzer lower bound.\n\n\n\n\n\\begin{proof}[Proof of Theorem~\\ref{theo:FR-LB}]\n\tPut $\\a = 2^{-73}$.\n\n\tWe have\n\t\\begin{equation}\\label{eq:FR-LB-ineq}\n\tc = 2^{-400} \\le \\a^4\/1500 \\;.\n\t\\end{equation}\n\n\n\tGiven $s$, let $H$ and $\\V_0$ be as in Theorem~\\ref{theo:main}.\n\tLet $\\P=(\\Z,\\E)$ be an $\\e$-regular $(\\ell,t,\\e_2(\\ell))$-equipartition of $H$\n\twith $\\e \\le c p$, $\\e_2(\\ell) \\le c \\ell^{-3}$ and $\\Z \\prec \\V_0$.\n\tThus, the bound $|Z| \\ge \\wow(s)$ would follow from Theorem~\\ref{theo:main}\n\n\tif we show that $\\P$ is an $\\langle \\a \\rangle$-regular partition of $H$.\n\n\n\tFirst we need to show that $\\P$ is $\\langle \\a \\rangle$-good (recall Definition~\\ref{def:k-good}). Let $E$ be a graph with $E \\in \\E$ on vertex classes $(Z,Z')$ (so $Z \\neq Z' \\in \\Z$).\n\tWe need to show that $E$ is $\\langle \\a \\rangle$-regular.\n\n\tSince $\\P$ is an $(\\ell,t,\\e_2(\\ell))$-equipartition we have\n\t(recall Definition~\\ref{def:ve-partition}) that $E$ is $\\e_2(\\ell)$-regular and $d(E) \\ge \\ell^{-1} - \\e_2(\\ell)$.\n\tThe statement's assumption on $\\e_2(\\ell)$ thus implies\n\t$d(E) \\ge 2\\e_2(\\ell)$.\n\tIt follows that for every $S \\sub Z$, $S' \\sub Z'$ with $|S| \\ge \\e_2(\\ell)|Z|$, $|S'| \\ge \\e_2(\\ell)|Z'|$ we have $d_E(S,S') \\ge d(E)-\\e_2(\\ell) \\ge \\frac12 d(E)$.\n\tThis proves that $E$ is $\\langle \\e_2(\\ell) \\rangle$-regular, and since $\\e_2(\\ell) \\le c \\le \\a$, that $E$ is $\\langle \\a \\rangle$-regular, as needed.\n\n\t\n\t\n\n\n\n\n\n\n\tIt remains to show that the $\\langle \\a \\rangle$-good $\\P$ is an $\\langle \\a \\rangle$-regular partition of $H$ (recall Definition~\\ref{def:k-reg}).\n\tBy symmetry, it suffices to show that $\\E_3 \\cup \\Z_3$ is an $\\langle \\a \\rangle$-regular partition of $G_{H}^3$.\n\tLet $H'$ be obtained from $H$ by removing all ($3$-)edges in triads of $\\P$ that are either not $\\e$-regular in $H$ or have density at most $3\\e$ in $H$.\n\tBy Definition~\\ref{def:FR-reg}, the number of edges removed from $H$ to obtain $H'$ is at most\n\t\\begin{equation}\\label{eq:FR-LB-modify}\n\t\\e|V(H)|^3 + 3\\e|V(H)|^3 \\le 4\\cdot c p |V(H)|^3\n\t\\le (\\a p\/27)|V(H)|^3 = \\a\\cdot e(H) \\;,\n\t\\end{equation}\n\n\twhere the second inequality uses ~(\\ref{eq:FR-LB-ineq}),\n\tand the equality uses the fact that all three vertex classes of $H$ are of the same size.\n\n\n\tThus, in $H'$, every non-empty triad of $\\P$ is $\\e$-regular and of density at least $3\\e$.\n\tPut $\\d = (\\a\/2)^2$.\n\tAgain by Definition~\\ref{def:FR-reg}, for every triad $P$ of $\\P$ and every subtraid $P'$ of $P$ with $t(P') \\ge \\d \\cdot t(P)$ ($\\ge \\e \\cdot t(P)$ by~(\\ref{eq:FR-LB-ineq})) we have $d_{H'}(P') \\ge d_{H'}(P)-\\e \\ge \\frac23 d_{H'}(P)$.\n\n\tIt follows from applying Claim~\\ref{claim:reduction} with $H'$ and $\\d$,\n\n\tusing~(\\ref{eq:FR-LB-ineq}),\n\n\n\tthat $\\E_3 \\cup \\Z_3$ is a perfectly $\\langle \\a \\rangle$-regular partition of $G_{H'}^3$.\n\tNote that~(\\ref{eq:FR-LB-modify}) implies that\n\tone can add\/remove at most $\\a \\cdot e(G_{H}^3)$ edges of $G_{H}^3$ to obtain $G_{H'}^3$.\n\tThus, $\\E_3 \\cup \\Z_3$ is an $\\langle \\a \\rangle$-regular partition of $G_{H}^3$, and as explained above, this completes the proof.\n\n\n\\end{proof}\n\n\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\nDengue Fever is a mosquito-borne disease with a significant global burden. Half of the world's population and 129 countries are at risk of infection \\cite{brady2012refining}. Between 100 and 400 million infections are registered each year globally \\cite{bhatt2013global}. Dengue presents as a flu-like illness with symptoms ranging from mild to severe. There are four serotypes, meaning up to four infections are possible for each person over their lifetime, with potentially life-threatening complications arising from severe Dengue. There is no specific treatment or universal vaccine, and the primary preventative measure is vector control. This underscores the importance of disease surveillance. Predictive models can help efficiently allocate public health resources to combat Dengue, with the goal of reducing the overall disease burden \\cite{who2020dengue}. Dengue incidence is associated with many different risk factors. These include climate conditions such as rainfall, extreme weather events, and temperature \\cite{chien2014impact, thu1998effect, fan2013identifying}, land use \\cite{kilpatrick2012drivers}, and poverty \\cite{mulligan2014dengue}. These factors may have nonlinear, context-specific, and time-variant effects on disease incidence, which poses a challenge to disease modeling. \n\nThe literature on Dengue forecasting is multi-disciplinary, uniting expertise from areas such as epidemiology, environmental science, computer science, and mathematics. Modeling frameworks include both theoretical and data-driven approaches. Compartmental models, for example, estimate the dynamics of Dengue infections in a population over time, and are based on extensive knowledge of the host, vector, and transmission process. Another simulation technique, agent based modeling is useful for estimating impacts of interventions, such as the release of sterile males to control the mosquito population \\cite{isidoro2009agent}. Statistical time series forecasting takes a more data-centric approach and is effective at modeling the highly auto-correlated nature of Dengue Fever \\cite{johansson2016evaluating, riley2020sarima_preprint}. Machine learning models leverage the increasing data availability on risk factors of disease spreading and offer non-parametric approaches that require less detailed knowledge of the disease and context \\cite{guo2017developing}. \n\nNeural networks are a subset of machine learning algorithms, which have made significant contributions to medicine and public health, including applications such as medical image analysis for disease diagnosis \\cite{qayyum2018medical}, identifying abnormalities in signals such as electrocardiographs (ECG) \\cite{papik1998application}, and optimizing decisions of health care providers, hospitals, and policy-makers \\cite{shahid2019applications}. \nNeural networks are also used to forecast diseases, including Malaria \\cite{zinszer2012scoping}, Influenza. \\cite{alessa2017review}, and Covid-19 \\cite{bullock2020mapping}. \n\nThis review examines the use of neural networks for Dengue Fever prediction. The objective is to summarize the existing literature and also provide an introduction to this still somewhat novel modeling technique. Our contributions are as follows:\n\\begin{itemize}\n\n \\item We summarize the technical decisions made in the literature, including architecture selection and hyper-parameter tuning. \n \\item We examine the data inputs (such as climate or population demographics) and best predictors of Dengue fever identified in specific contexts.\n \\item We review the relative performance of different neural network architectures and comparator models, such as other machine learning techniques. \n\\end{itemize}\n\n\n\nTo the knowledge of the authors, no systematic review of the literature on neural networks applied to Dengue Fever prediction has yet been conducted. \nSiriyasatien et al (2018) \\cite{siriyasatien2018dengue} provide a broader review of data science models applied for Dengue Fever prediction. Racloz et al (2012) \\cite{racloz2012surveillance} review the literature on surveillance systems for Dengue Fever. The study finds that most papers use logistic or multiple regression to analyze Dengue risk factors. (Seasonal) Auto-Regressive Integrated Moving Average (S\/ARIMA) models have become a popular choice to incorporate auto-regressions but are not suitable for all data types. Though also apt at analyzing auto-regressive behavior, the study does not review any literature on neural networks or machine learning for Dengue surveillance. We aim to fill this knowledge gap in the present review.\n\n\n\n\n\n\\section{A Primer on Neural Networks}\n\nThis section gives an overview of some of the neural network models that are relevant to Dengue prediction. We explain the model first intuitively, then mathematically, and focus on feed-forward neural networks, which are used most often for Dengue prediction. Many high-quality textbooks and open source tutorials offer a deeper introduction to neural networks (for example \\cite{deng2014deep, haykin1999neural}). Several Python libraries include easy-to-use implementations of neural network models (for example TensorFlow \\cite{tensorflow}, scikit-learn \\cite{scikit-learn}, or Keras \\cite{chollet2015keras}). We also provide sample code for Dengue prediction in a Github repository\\footnote{https:\/\/github.com\/KRoster\/NN4Dengue}.\n\n\nLike other machine learning models, neural networks learn to execute specific supervised tasks, such as predicting the number of Dengue infections, based on a large set of labeled examples. Given a set of inputs, such as climate conditions, the model learns to estimate the output, such as next month's Dengue incidence. After making a first guess, the model looks at the correct answer and updates its parameters to reduce its error on the next iteration until its predictions are optimized.\n\nComputationally, neural networks are represented as a network of processing units (\"neurons\" or \"hidden units\") that transform and pass information from one layer to the next. During training, we distinguish between forward- and backward-propagation, depending on the direction of the information flow. During the forward-propagation step, input information passes between units, each time being transformed according to a set of weights and a non-linear activation function. The prediction error of the output relative to the true label is computed. Back-propagation is then used to reduce this error on the next iteration: The weights of each unit, the parameters that define how information is combined, are updated according to their influence on the prediction error. This combination of forward- and backward-propagation is repeated several times until the predictions are increasingly accurate.\n\nAfter this training phase, model performance is tested on a hold-out test set. Since the test set is not used to train the model, it can give a good indication of how well the model will generalize to new data, once applied in the real world.\n\n\\begin{figure}[!h]\n \\centering\n \\includegraphics[width=0.9\\textwidth]{\"nn_fully_annotated\".png}\n \\caption{Sample feed-forward neural network architecture}\n \\medskip\n \\small\n The figure shows a sample neural network architecture with 4 input features, 2 hidden layers with 5 and 2 nodes respectively and a single node in the output layer.\n \\label{fig:nn annotated}\n\\end{figure}\n\nThere are many different neural network architectures that are designed for different kinds of inputs and tasks. The architecture determines how the units are connected, how information flows through the network, and how many weights need to be optimized. An example of a feed-forward network is illustrated in figure \\ref{fig:nn annotated}. It shows the forward propagation of information through the individual neurons of the network, which was described above. The information that is passed through a given neuron $j$ in layer $l$, called its activation $a_{j}^{[l]}$, is computed as the linear combination of weights $w_{ij}^{[l]}$, biases $b_{j}^{[l]}$, and inputs from the previous layer $a_{i}^{[l-1]}$, which is passed through a nonlinear activation function $g^{[l]}$: \n\n\\begin{flalign} \\label{eq forward prop}\n a_{j}^{[l]} &= g^{[l]}(\\sum\\limits_{i=1}^{n_l} w_{ij}^{[l]} a_{i}^{[l-1]} + b_{j}^{[l]}) \\\\\n \\text{with: } a_{j}^{[0]} &= X_{j} \\\\\n \\hat{y} &= a^{[k]}\n\\end{flalign}\n\n\\noindent where ${n_l}$ is the number of neurons in layer $l$, $k$ is the number of layers (the input layer is counted as layer $0$), and $\\hat{y}$ is the prediction that is generated in the final layer.\n\n\nThe gradient of the loss with respect to a given weight $w_{ij}^{[l]}$ or bias $b_{j}^{[l]}$ tells the model how the given parameter needs to be adjusted, with the learning rate $\\alpha$ determining the size of the adjustment:\n\n\\begin{flalign} \\label{eq gradient descent}\n \\theta &= \\theta - \\alpha \\frac{d\\mathcal{L}}{d\\theta}\n\\end{flalign}\n\n\\noindent where $\\theta$ is a given set of weights or biases to be updated and $\\mathcal{L}$ is the loss.\n\n\n\nA benefit of neural networks is that parameters are updated by the model itself, not predefined by the researcher. They do not rely on strong assumptions and knowledge of the disease. Yet some hyper-parameters of neural network models must still be set by the researcher and are generally determined through a combination of domain knowledge and iterative experimentation. These hyper-parameters are tuned using a validation set or through cross-validation. Hyper-parameters may include the learning rate (by how much the weights are updated at each backward propagation step), the number of epochs (how often the forward and backward propagation steps are repeated), the mini-batch size (how many training examples are processed at each step), how weights are initialized (for example all zeros or random values), the number of hidden units in each layer, and the number of layers. Other important decisions include the choice of activation function in each layer, the choice of loss function, and the relative size of training, validation, and test sets. The optimal choices and ranges tested for these hyper-parameters depend, among other aspects, on the size and type of data available and the nature of the predictive task. One aim of this review is to summarize the choice of parameters deemed optimal by researchers in the existing literature to assist future researchers in the hyper-parameter tuning process. \n\nBesides feed-forward neural networks, this review includes two other categories of models used for Dengue forecasting. Recurrent neural networks (RNN) were developed for the analysis of sequence data, such as time series (sequences of observations in time) or text (sequences of words). They are cyclical (they contain loops) and process sequences of inputs together, giving the network a \"memory\" to remember historical feature states. RNNs are used in contexts where the evolution of input features matters for the target feature, where the prediction relies on information at multiple time steps. An example is machine translation, where the meaning of a word is influenced by its context, by the other words in the sentence. The number of Dengue cases may also depend on the sequence of risk factors or the sequence of previous Dengue incidence. Using time series data, RNNs may help capture the auto-correlation of the disease and lagged influence of risk factors, such as rainfall or travel patterns. Within the family of RNNs we distinguish different architectures, such as Long Short-Term Memory (LSTM) and Gated Recurrent Units (GRU) \\cite{hochreiter1997lstm, schmidhuber2015deep, lipton2015critical}. \n\nConvolutional neural networks (CNN) are another subclass of deep neural network models with applications in computer vision \\cite{lecun1990handwritten}, such as image segmentation and classification. Inputs into CNNs are three-dimensional - such as a two dimensional image with three color channels, which form a 3-D matrix of pixel values. The architectures of CNNs include different kinds of operations, such as convolution and pooling, which help identify key regions in the image that help make the final prediction. For example, the model may learn that the presence of a body of water (shiny pixels in a satellite image) influence the number of Dengue cases we should expect. Crucial in this example is that the type of relevant landscape feature is not defined by the researcher but learned by the model. However, CNNs may also be used for data processing prior to prediction, for example to classify land use (such as vegetation or urban areas) from satellite images, which are known to influence Dengue spread and which can then be fed into Dengue a separate predictive model.\n\n\nThe tasks executed by machine learning models in general, and neural networks in particular, fall into two different categories: classification and prediction. Classification is the process of assigning a categorical label to an example based on its inputs, such as determining whether a patient has Dengue fever or a different disease. Prediction is the process of forecasting a future state based on historical or current inputs, for example predicting next month's Dengue incidence from this month's rainfall. \n\nThis review is focused on predictive models of Dengue Fever at the population level. In addition to classic regression models, we include papers using classification for prediction approaches, which produce categorical instead of continuous outputs. For example, a model may predict whether a city will experience an outbreak (binary classification) or which risk category will prevail (multi-label classification). \n\n\n\n\n\n\n\n\n\\section{Methodology}\n\n\\begin{figure}[!h]\n \\centering\n \\includegraphics[width=0.8\\textwidth]{\"PRISMA_2009Flow_filled_new\".pdf}\n \\caption{PRISMA Flow Chart}\n \\label{fig:prisma flow}\n\\end{figure}\n\nThis review follows the Preferred Reporting Items for Systematic Reviews and Meta-Analyses (PRISMA) statement guidelines \\cite{prisma2009}. The PRISMA flow chart is represented in figure \\ref{fig:prisma flow}. A systematic search was conducted using Web of Science\/ Knowledge, Scopus (abstract, title, keywords), PubMed, and Science Direct (abstract, title) databases. References of papers appearing in the search results were also examined for relevant works. The searches were conducted between 21 July and 17 August 2020 and used the following search string:\n\\begin{center}\n \"( deep learning OR neural network ) AND dengue\"\n\\end{center}\n\n\\begin{comment}\nThe databases yielded the following numbers of results:\n\n\\begin{itemize}\n \\item Web of Science \/ Knowledge: 92 (abstract screened: 46) (full text review: 26)\n \\item PubMed: 21 (abstract screened: 4) (full text review: 4)\n \\item Scopus: 23 (abstract screened: 18) (full text review: 8)\n \\item Science Direct: 13 (abstract screened: 3) (full text review: 3)\n \\item references of other papers: 31 (abstract screened: 27) (full text review: 21)\n\\end{itemize}\n\n\\end{comment}\n\n\n\n\n\n\nThe papers were examined for inclusion in the review in three stages: by title, by abstract, and finally by full text. In each phase, the following inclusion and exclusion criteria were applied liberally to determine which papers would be considered in the next phase. If the information in the title and\/or abstract was inconclusive, the paper was included in the full-text review. Inclusion and exclusion criteria are as follows:\n\n\\begin{itemize}\n \\item Studies must implement a neural network or deep learning technique, either as the main method or as a comparator model. Reviews of the literature on neural networks for Dengue forecasting would also be considered. \n \\item Studies must predict Dengue fever incidence or risk. There are no restrictions on the type of target variable used. For example, the number of Dengue cases, Dengue incidence rate, or a binary Dengue risk variable are all accepted as target features.\n \\item Studies must examine Dengue in a human population. Models for disease diagnosis of individuals are excluded. Studies modeling the location of vectors without relation to Dengue incidence are excluded. Studies examining animal hosts are excluded.\n \\item As the search query is in English, only English language articles were identified and included in the review.\n\\end{itemize}\n\n\n\n\n\n\\section{Results}\nThis section summarizes the findings of the review. The full list of papers and their properties are provided in table 1 at the end of this section.\n\n\n\\subsection{Prediction Target}\nWithin the criteria of study selection, there was significant variation in the study specifications, including the formulation of the target feature. Most studies predict the future number of Dengue cases based on historical time series. However, there are some notable exceptions. Andersson et al (2019) \\cite{andersson2019combining} predict a static incidence rate at the neighborhood level in Rio de Janeiro, Brazil. Two studies predict the risk of a Dengue outbreak: Abeyrathna et al (2019) \\cite{abeyrathna2019scheme} formulate their task as a binary classification of 'outbreak' or 'no outbreak', while Anno et al (2019) \\cite{anno2019spatiotemporal} classify cities according to five risk categories. \n\nEarly forecasts provide more time for public health response but may have lower confidence. In most of the selected studies, researchers predict one period ahead corresponding to the measurement frequency of the data used. A single study predicted the present: Livelo and Cheng (2018) \\cite{livelo2018intelligent} use social media activity in the Philippines to predict the present situation of Dengue infection. They use a neural network model to classify tweets according to a set of different categories of Dengue-related topics. Their weekly Dengue index based on this Twitter data is correlated with actual Dengue case counts. Koh et al. (2018) \\cite{koh2018model} compare one and two-week forecast horizons. Soemsap et al. (2014) \\cite{soemsap2014forecasting} use a two-week horizon with weekly data inputs. Dharmawardhana et al (2017) \\cite{dharmawardana2017predictive} predict Dengue cases in Sri Lankan districts four weeks ahead using weekly data on cases, climate, human mobility, and vegetation indices. \nChakraborty et al. (2019) \\cite{Chakraborty2019forecasting} use the longest forecasting horizon. They use three different datasets of Dengue cases, two of which are weekly (Peru and Puerto Rico), and one which is measured monthly (Philippines). For the weekly data, they compare three and six month horizons. For the monthly data, they compare a horizon of six months with one year forecasts. For Peru and the Philippines, the best model changes with the prediction horizon. \n\n\n\n\n\\subsection{Data Sources}\n\n\\begin{figure}[!h]\n \\centering\n \\includegraphics[width=0.9\\textwidth]{\"articles_for_review_variables_used2\".pdf}\n \\caption{Variable types used in the selected studies}\n \\medskip\n \\small\n The figure illustrates how frequently studies included individual data types as well as combinations thereof. Each row represents a data type and the horizontal bars measure their frequency of inclusion in the selected papers. The columns represent the combinations of variables, and the vertical bars measure how often the given combinations occurred. \n \\label{fig:variables}\n\\end{figure}\n\nBesides epidemiological data on the history of Dengue Fever, which is used in all studies, the most common predictors are meteorological variables (used in 14 studies) (see figure \\ref{fig:variables}). Additional data sources include: vegetation indices and other landscape features extracted from satellite images (in two studies), human mobility, specifically mobile phone data (used in one study), aerial and street view images (in one study), social media data, specifically from Twitter (one study), vector information (one study), and demography (one study). Three papers use only epidemiological data.\n\n\n\\begin{figure}[!h]\n \\centering\n \\includegraphics[width=0.8\\textwidth]{\"articles_for_review_distribution_of_geographic_scale\".png}\n \\caption{Geographic granularity used in the selected studies}\n \\label{fig:geographic granularity}\n\\end{figure}\n\nMost studies examine granular geographic regions, such as cities (nine studies) or districts (six studies). Two studies even modeled Dengue at the sub-city (neighborhood) level (see figure \\ref{fig:geographic granularity}). Granular scales can be beneficial in two ways: they allow for local predictions and accordingly targeted public health response. They also generally result in larger training datasets with more observations than national or state-level aggregations, which plays to the strength of neural networks. The studies cover countries in both Asia and Latin America, which are considered the primary Dengue endemic regions \\cite{who2020dengue}. A total of 13 countries are included (see figure \\ref{fig:countries map}), with the Philippines and Brazil appearing most often. \n\n\n\n\\begin{figure}[!h]\n \\centering\n \\includegraphics[width=0.9\\textwidth]{\"articles_for_review_geographic_distribution\".png}\n \\caption{Countries covered by the selected studies}\n \\label{fig:countries map}\n\\end{figure}\n\n\n\n\n\\subsection{Model Selection}\n\n\\begin{figure}[!h]\n \\centering\n \\includegraphics[width=0.9\\textwidth]{\"articles_for_review_models_used2\".pdf}\n \\caption{Models used in the selected studies}\n \\label{fig:models}\n \\medskip\n \\small\n The figure illustrates how frequently studies used different forecasting models as well as combinations thereof. Each row represents an algorithm and the horizontal bars measure their frequency of inclusion in the selected papers. Green rows refer to neural network models, while blue bars show comparator models, such as machine learning algorithms. The columns illustrate which combinations of models were used, the frequency of which is measured by the vertical bars. \n\\end{figure}\n\n\nMost studies implement a shallow neural network with just one hidden layer. Four studies use deeper feed-forward architectures. Three studies employ recurrent neural networks, two of which used LSTM units and one used a GRU. Three studies implemented CNNs. \n\n\nThe neural network models were compared to a range of other models. Support Vector Machine (SVM) models were most common (six studies). Generalized Additive Models (GAMs), (S)ARIMA and linear autoregressive (AR) models occurred twice each. The remaining comparators were used in just one paper each: Gradient Boosting Machine (GBM), XGBoost, Linear regression, Poisson regression, Tsetlin Machine (TM), Cellular automata, a compartmental model, and a naive baseline. Figure \\ref{fig:models} shows how frequently different model types and combinations thereof were used in the selected literature. \n\nThe most common comparator model used is SVM. Evidence suggests that both SVM and neural networks perform well, and which is better depends on the context. In Xu et al (2020) \\cite{xu2020forecast}, an LSTM neural network outperforms an SVM model in predicting the number of Dengue fever cases in 20 Chinese cities. Transfer learning further increased the performance gains in cities with few overall Dengue cases. Abeyrathna et al (2019) \\cite{abeyrathna2019scheme} use a classification for prediction approach to identify whether districts in the Philippines are likely to experience a Dengue Fever outbreak. Whether SVM or the neural network model is better in the context of this study depends on the evaluation metric used: SVM has a higher precision, while ANN has a higher F1-score. Yusof and Mustaffa (2011) \\cite{yusof2011dengue} and Kesorn et al (2015) \\cite{kesorn2015morbidity} both find SVM to have a higher performance than the neural network. \n\nFour studies use a feed-forward network with at least two hidden layers. They provide varying evidence to the performance relative to shallow neural networks or other machine learning models. In \\cite{dharmawardana2017predictive}, a model with two hidden layers performs better than an XGBoost model in forecasting Dengue in Sri Lankan districts. In two studies (\\cite{abeyrathna2019scheme} and \\cite{baquero2018dengue}), the deep neural network provides better predictions than alternative neural network models, but is not the best model overall: In \\cite{baquero2018dengue}, the MLP with two hidden layers performs better than the LSTM-RNN to forecast Dengue in a Brazilian city, but a GAM is best overall. In \\cite{abeyrathna2019scheme}, a model with three hidden layers has lower errors than models with a single or five hidden layers. However, the best F1 score in classifying binary outbreak risk is achieved by a Tsetlin Machine. The final study \\cite{rahayu2019prediction} uses a model with three hidden layers but does not employ any comparator models.\n\n\nPerformance of the RNN models is also mixed. The two studies using LSTM-RNNs produce contradicting results despite having similar study designs. Baquero et al. (2018) \\cite{baquero2018dengue} and Xu et al. (2020) \\cite{xu2020forecast} both use LSTM-RNNs to forecast Dengue at the city level, compare their model to a feed-forward neural network and a GAM, use meteorological and epidemiological input features, and evaluate performance using RMSE. Xu et al. (2020) \\cite{xu2020forecast} find the LSTM to have the best performance relative to all comparator models, whereas in the study by Baquero et al. (2018) \\cite{baquero2018dengue}, the LSTM-ANN performs worse than the simple neural network and the best model overall is the GAM. \n\nThe three studies using CNNs do not provide comparisons to other types of models, though one study compares different CNN architectures. However, the key advantage of the examined CNN studies is the granular geographic resolution in predictions. The models take satellite and other images as inputs and can therefore generate forecasts at the sub-city level. Andersson et al. (2019) \\cite{andersson2019combining} compare different kinds of CNN architectures (based on DenseNet-161) for each of three input types, specifically satellite imagery, street view images, and the combination of the two. Their model predicts Dengue incidence rates at the neighborhood level in Rio de Janeiro, Brazil. Best performance is achieved when combining street and aerial images. Rehman et al. (2019) \\cite{rehman2019deep} implement a pre-trained CNN (based on U-Net architecture) to extract landscape features from satellite imagery, which they use in a compartmental SIR model to forecast Dengue in neighborhoods of two Pakistani cities. Anno et al. (2019) \\cite{anno2019spatiotemporal} use a CNN (based on AlexNet architecture) to classify five levels of outbreak risk in cities in Taiwan, based on images capturing sea surface temperatures. \n\n\n\n\n\n\n\n\\subsection{Architecture}\nArchitectures vary across the different models. The studies implementing shallow models with one hidden layer cumulatively tested 1-30 hidden units. Three studies (\\cite{wijekoon2014prediction, aburas2010dengue, soemsap2014forecasting}) tested a broad number of values across this range, and their optimal number of hidden units were related to the number of input features. Two studies \\cite{wijekoon2014prediction, aburas2010dengue} identified four hidden units as the optimum. Both used meteorological and epidemiological variables and both measured four different pieces of information. Wijekoon et al (2014) \\cite{wijekoon2014prediction} used temperature, rainfall, humidity and Dengue cases. Aburas et al (2010) \\cite{aburas2010dengue} included a total of seven features that measured the same information but with additional lags. The third study \\cite{soemsap2014forecasting} identified 25 nodes as the optimum, which matches the larger input size of 32 features. Instead of trying different values, Datoc et al (2016) \\cite{datoc2016forecasting} fixed the number of hidden units to the number of input features, and selected the best-performing model among eight models with three to five nodes. The best model used three nodes (corresponding to three input features). \n\nAbeyrathna et al (2019) \\cite{abeyrathna2019scheme} compared shallow and deep architectures. Their best neural network model had three layers with (20, 150, 100) units respectively. This model performed better than the two shallow models with five or 20 neurons in a single hidden layer. It also outperformed a five-hidden layer model of (20, 200, 150, 100, 50) topology. \n\nBaquero et al (2018) \\cite{baquero2018dengue} also used deep architectures with 2 layers with combinations of (10 or 20) neurons in the first layer, and (5 or 10) neurons in the second. They implemented both MLP and LSTM models. The best neural network model was an MLP with (20,10) neurons, though a Generalized Additive Model (GAM) and an ensemble both performed better than the neural networks. \n\nThe other LSTM model was implemented by Xu et al (2020) \\cite{xu2020forecast} with a different architecture. They used a single hidden layer with 64 neurons. Unlike Baquero et al (2018) \\cite{baquero2018dengue}, their model outperformed GAM as well as SVR and GBM. \n\n\n\n\\subsection{Model Evaluation}\nThe studies utilize a range of metrics to assess model performance, in part due to different prediction targets. The most common measures are Root Mean Squared Error (RMSE) and Pearson's correlation, which are each used in seven studies each. Mean Average Error (MAE), Mean Squared Error (MSE), and accuracy are also commonly used.\n\n\n\n\n\n\\section{Discussion}\n\nThe literature on neural networks applied to Dengue forecasting is still somewhat scarce, but the reviewed studies include some promising examples, where neural networks outperform other machine learning or statistical approaches. The literature therefore suggests that neural networks may be appropriate for inclusion in a set of candidate models when forecasting Dengue incidence. However there is variation as to whether neural network models or other approaches perform best. Future research may compare different model types and architectures on multiple datasets, to better understand how the ideal model varies by context, such as geography or data availability.\n\nThough many risk factors of Dengue have been identified, most neural network models limit inputs to meteorological and epidemiological data. Future studies may evaluate the value of alternative predictors, especially as NNs tend to deal well with high-dimensional problems. Along the same lines, there is space for more incorporation of non-causal predictors of Dengue fever from sources such as social media. Some papers in this review have leveraged these new predictors, specifically from Twitter \\cite{livelo2018intelligent}, Street view and satellite images \\cite{andersson2019combining}, mobile phone records \\cite{dharmawardana2017predictive}, and features derived from satellite imagery \\cite{dharmawardana2017predictive, rehman2019deep}. These data sources have also been effective with other machine learning models, for example Baidu search data to predict Dengue in a Chinese province \\cite{guo2017developing}. \n\nTransfer learning was implemented in \\cite{xu2020forecast}, where it provided promising results. The authors trained a model on a city with high Dengue incidence and used it to predict disease in lower-incidence geographies. This was the only study in this review that used transfer learning outside the context of pre-trained CNNs. This may present another avenue for further research, especially in locations where data is scarce. \n\n\\subsubsection*{Acknowledgments}\n\nKirstin Roster gratefully acknowledges support from the S\u00e3o Paulo Research Foundation (FAPESP) (grant number 2019\/26595-7). Francisco Rodrigues acknowledges partial support from CNPq (grant number 309266\/2019-0) and FAPESP (grant number 2019\/23293-0).\n\n\\begin{comment}\n\\begin{table}[]\n \\centering\n \n \n {\\small %\n \\begin{tabular}{L{3cm} | L{1.5cm} | L{1.5cm} | L{2cm} | L{1.5cm} | L{3cm} }\n \\toprule\n Reference & Country & Geospatial scale & Target & Prediction horizon & Input data \\\\\n \\midrule\n \n \\bottomrule\n \\end{tabular}\n }%\n \n\\end{table}\n\\end{comment}\n\n\n\\includepdf[pages={1-}, fitpaper]{\"Tableofincludedpapers_2021-05-31\".pdf}\n\n\n\\newpage\n\\nocite{*}\n\\bibliographystyle{unsrt}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nSubmodular set functions are defined by the following condition for all pairs of sets $S,T$:\n$$ f(S \\cup T) + f(S \\cap T) \\leq f(S) + f(T),$$\nor equivalently by the property that the {\\em marginal value} of any element,\n$f_S(j) = f(S+j)-f(S)$, satisfies $f_T(j) \\leq f_S(j)$, whenever\n$j \\notin T \\supset S$. In addition, a set function is called monotone\nif $f(S) \\leq f(T)$ whenever $S \\subseteq T$.\nThroughout this paper, we assume that $f(S)$ is nonnegative.\n\nSubmodular functions have been studied in the context of combinatorial\noptimization since the 1970's, especially in connection with matroids\n\\cite{E70,E71,NWF78,NWF78II,NW78,W82a,W82b,L83,F97}.\nSubmodular functions appear mostly for the following two reasons:\n(i) submodularity arises naturally in various combinatorial settings,\nand many algorithmic applications use it either explicitly or implicitly;\n(ii) submodularity has a natural interpretation as the property\nof {\\em diminishing returns}, which defines an important class of\nutility\/valuation functions.\nSubmodularity as an abstract concept is both general enough to be useful for applications\nand it carries enough structure to allow strong positive results.\nA fundamental algorithmic result is that any submodular function\ncan be {\\em minimized} in strongly polynomial time \\cite{FFI01,Lex00}.\n\nIn contrast to submodular minimization, submodular maximization problems are\ntypically hard to solve exactly. Consider the classical problem\nof maximizing a monotone submodular function subject to a cardinality constraint,\n$\\max \\{f(S): |S|\\leq k\\}$. It is known that this problem admits a $(1-1\/e)$-approximation\n\\cite{NWF78} and this is optimal for two different reasons:\n(i) Given only black-box access to $f(S)$, we cannot achieve a better approximation,\nunless we ask exponentially many value queries \\cite{NW78}.\nThis holds even if we allow unlimited computational power.\n(ii) In certain special cases where $f(S)$ has a compact representation on the input,\nit is NP-hard to achieve an approximation better than $1-1\/e$ \\cite{Feige98}.\nThe reason why the hardness threshold is the same in both cases\nis apparently not well understood.\n\nAn optimal $(1-1\/e)$-approximation for the problem $\\max \\{f(S): |S| \\leq k\\}$ where $f$ is\nmonotone sumodular is achieved\nby a simple greedy algorithm \\cite{NWF78}. This seems to be rather coincidental; for\nother variants of submodular maximization, such as unconstrained (nonmonotone)\nsubmodular maximization \\cite{FMV07}, monotone submodular maximization subject\nto a matroid constraint \\cite{NWF78II,CCPV07,Vondrak08}, or submodular\nmaximization subject to linear constraints \\cite{KTS09,LMNS09},\ngreedy algorithms achieve suboptimal results. A tool which has proven useful\nin approaching these problems is the {\\em multilinear relaxation}.\n\n\\\n\n\\noindent{\\bf Multilinear relaxation.}\nLet us consider a discrete optimization problem $\\max \\{f(S): S \\in {\\cal F}\\}$, where\n$f:2^X \\rightarrow {\\boldmath R}$ is the objective function and ${\\cal F} \\subset 2^X$\nis the collection of feasible solutions. In case $f$ is a linear function,\n$f(S) = \\sum_{j \\in S} w_j$, it is natural to replace this problem by a linear\nprogramming problem. For a general set function $f(S)$, however, a linear\nrelaxation is not readily available. Instead, the following relaxation\nhas been proposed \\cite{CCPV07,Vondrak08,CCPV09}:\nFor ${\\bf x} \\in [0,1]^X$, let $\\hat{{\\bf x}}$ denote a random vector in $\\{0,1\\}^X$\nwhere each coordinate of $x_i$ is rounded independently to $1$ with probability $x_i$\nand $0$ otherwise.\\footnote{We denote vectors consistently in boldface $({\\bf x})$ and their coordinates in italics $(x_i)$.\nWe also identify vectors in $\\{0,1\\}^n$ with subsets of $[n]$ in a natural way.\n}\nWe define\n$$ F({\\bf x}) = {\\bf E}[f(\\hat{{\\bf x}})] = \\sum_{S \\subseteq X} f(S) \\prod_{i \\in S} x_i\n \\prod_{j \\notin S} (1-x_j).$$\nThis is the unique {\\em multilinear polynomial} which coincides with $f$ on $\\{0,1\\}$-vectors.\nWe remark that although we cannot compute the exact value of $F({\\bf x})$ for a given\n${\\bf x} \\in [0,1]^X$ (which would require querying $2^n$ possible values of $f(S)$),\nwe can compute $F({\\bf x})$ approximately, by random sampling. Sometimes this causes\ntechnical issues, which we also deal with in this paper.\n\nInstead of the discrete problem $\\max \\{f(S): S \\in {\\cal F}\\}$,\nwe consider a continuous optimization problem $\\max \\{F({\\bf x}): {\\bf x} \\in P({\\cal F})\\}$,\nwhere $P({\\cal F})$ is the convex hull of characteristic vectors corresponding to ${\\cal F}$,\n$$ P({\\cal F})\n = \\left\\{ \\sum_{S \\in {\\cal F}} \\alpha_S {\\bf 1}_S:\n \\sum_{S \\in {\\cal F}} \\alpha_S = 1, \\alpha_S \\geq 0 \\right\\}.$$\nThe reason why the extension $F({\\bf x}) = {\\bf E}[f(\\hat{{\\bf x}})]$ is useful\nfor submodular maximization problems is that $F({\\bf x})$ has convexity properties\nthat allow one to solve the continuous problem $\\max \\{F({\\bf x}): {\\bf x} \\in P\\}$\n(within a constant factor) in a number of interesting cases.\nMoreover, fractional solutions can be often\nrounded to discrete ones without losing {\\em anything} in terms of the objective\nvalue. Then, our ability to solve the multilinear relaxation\ntranslates directly into an algorithm for the original problem.\nIn particular, this is true when the collection of feasible solutions\nforms a matroid.\n\n{\\em Pipage rounding} was originally developed by Ageev and Sviridenko\nfor rounding solutions in the bipartite matching polytope \\cite{AS04}.\nThe technique was adapted to matroid polytopes by Calinescu et al.\n \\cite{CCPV07}, who proved that for any submodular function $f(S)$\nand any ${\\bf x}$ in the matroid base polytope $B({\\cal M})$,\nthe fractional solution ${\\bf x}$ can be rounded to a base\n$B \\in {\\cal B}$ such that $f(B) \\geq F({\\bf x})$.\nThis approach leads to an optimal $(1-1\/e)$-approximation for\nthe Submodular Welfare Problem,\nand more generally for monotone submodular maximization subject\nto a matroid constraint \\cite{CCPV07,Vondrak08}.\nIt is also known that\nthe factor of $1-1\/e$ is optimal for the Submodular Welfare Problem\nboth in the NP framework \\cite{KLMM05} and in the value oracle model\n\\cite{MSV08}. Under the assumption that the submodular function $f(S)$\nhas {\\em curvature} $c$, there is a $\\frac{1}{c}(1-e^{-c})$-approximation\nand this is also optimal in the value oracle model \\cite{VKyoto08}.\nThe framework of pipage rounding can be also extended to nonmonotone submodular functions;\nthis presents some additional issues which we discuss in this paper.\n\nFor the problem of\nunconstrained (nonmonotone) submodular maximization, a $2\/5$-approximation\nwas developed in \\cite{FMV07}. This algorithm implicitly uses\nthe multilinear relaxation $\\max \\{F({\\bf x}): {\\bf x} \\in [0,1]^X\\}$.\nFor {\\em symmetric} submodular functions, it is shown in \\cite{FMV07}\nthat a uniformly random solution ${\\bf x} = (1\/2,\\ldots,1\/2)$ gives\n$F({\\bf x}) \\geq \\frac12 OPT$, and there is no better approximation algorithm\nin the value oracle model. Recently, a $1\/2$-approximation was found\nfor unconstrained maximization of a general nonnegative submodular function \\cite{BFNS12}.\nThis algorithm can be formulated in the multilinear relaxation framework, but also\nas a randomized combinatorial algorithm.\n\nUsing additional techniques, the multilinear relaxation can be also applied\nto submodular maximization with knapsack constraints\n($\\sum_{j \\in S} c_{ij} \\leq 1$). For the problem of maximizing a monotone\nsubmodular function subject to a constant number of knapsack constraints,\nthere is a $(1-1\/e-\\epsilon)$-approximation algorithm for any $\\epsilon > 0$ \\cite{KTS09}.\nFor maximizing a nonmonotone submodular function\nsubject to a constant number of knapsack constraints,\na $(1\/5-\\epsilon)$-approximation was designed in \\cite{LMNS09}.\n\nOne should mention that not all the best known results for submodular maximization\nhave been achieved using the multilinear relaxation. The greedy algorithm\nyields a $1\/(k+1)$-approximation for monotone submodular maximization\nsubject to $k$ matroid constraints \\cite{NWF78II}. Local search methods\nhave been used to improve this to a $1\/(k+\\epsilon)$-approximation,\nand to obtain a $1\/(k+1+1\/(k-1)+\\epsilon)$-approximation\nfor the same problem with a nonmonotone submodular function,\nfor any $\\epsilon>0$ \\cite{LMNS09,LSV09}.\nFor the problem of maximizing a nonmonotone submodular function\nover the {\\em bases} of a given matroid, local search yields\na $(1\/6-\\epsilon)$-approximation, assuming that the matroid contains two disjoint bases\n\\cite{LMNS09}.\n\n\n\\subsection{Our results}\n\nOur main contribution (Theorem~\\ref{thm:general-hardness}) \nis a general hardness construction that yields\ninapproximability results in the value oracle model in an automated\nway, based on what we call the {\\em symmetry gap} for some fixed instance.\nIn this generic fashion, we are able to replicate a number of previously\nknown hardness results (such as the optimality of the factors\n$1-1\/e$ and $1\/2$ mentioned above),\nand we also produce new hardness results using this construction\n(Theorem~\\ref{thm:matroid-bases}).\nOur construction helps explain the particular hardness thresholds obtained under\nvarious constraints, by exhibiting a small instance where the threshold can be\nseen as the gap between the optimal solution and the best symmetric solution.\nThe query complexity results in \\cite{FMV07,MSV08,VKyoto08}\ncan be seen in hindsight as special cases of Theorem~\\ref{thm:general-hardness},\nbut the construction in this paper is somewhat different and\ntechnically more involved than the previous proofs for particular cases.\n\n\\paragraph{Concrete results}\nBefore we proceed to describe our general hardness result, we present its implications\nfor two more concrete problems.\nWe also provide closely matching approximation algorithms for these two problems,\nbased on the multilinear relaxation. In the following, we assume that the objective function is given by a value oracle\nand the feasibility constraint is given by a membership oracle: a value oracle for $f$ returns the value $f(S)$ for a given set $S$,\nand a membership oracle for ${\\cal F}$ answers whether $S \\in {\\cal F}$ for a given set $S$.\n\nFirst, we consider the problem of maximizing a nonnegative (possibly nonmonotone) submodular function\nsubject to a \\emph{matroid base} constraint. (This generalizes for example the maximum bisection\nproblem in graphs.) We show that the approximability of this problem is related to base\npackings in the matroid. We use the following definition.\n\n\\begin{definition}\n\\label{def:fractional-base}\nFor a matroid ${\\cal M}$ with a collection of bases ${\\cal B}$, the fractional base\npacking number is the maximum possible value of $\\sum_{B \\in {\\cal B}} \\alpha_B$\nfor $\\alpha_B \\geq 0$ such that $\\sum_{B \\in {\\cal B}: j \\in B} \\alpha_B \\leq 1$ for\nevery element $j$.\n\\end{definition}\n\n\\noindent{\\bf Example.}\nConsider a uniform matroid with bases ${\\cal B} = \\{ B \\subset [n]: |B| = k \\}$. ($[n]$ denotes the set of integers\n$\\{1,2,\\cdots,n\\}$.)\nHere, we can take each base with a coefficient $\\alpha_B = 1 \/ {n-1 \\choose k-1}$, which satisfies\nthe condition $\\sum_{B \\in {\\cal B}: j \\in B} \\alpha_B \\leq 1$ for every element $j$ since every element\nis contained in ${n-1 \\choose k-1}$ bases. We obtain that the fractional packing number is at least\n${n \\choose k} \/ {n-1 \\choose k-1} = \\frac{n}{k}$. It is also easy to check that the fractional packing number\ncannot be larger than $\\frac{n}{k}$.\n\n\\begin{theorem}\n\\label{thm:matroid-bases}\nFor any $\\nu$ in the form $\\nu = \\frac{k}{k-1}, k \\geq 2$,\nand any fixed $\\epsilon>0$, a $(1-\\frac{1}{\\nu}+\\epsilon) = \\frac{1}{k}$-approximation for the problem\n$\\max \\{f(S): S \\in {\\cal B}\\}$, where $f(S)$ is a nonnegative submodular function,\nand ${\\cal B}$ is a collection of bases in a matroid with fractional packing number at least $\\nu$,\nwould require exponentially many value queries.\n\nOn the other hand, for any $\\nu \\in (1,2]$,\nthere is a randomized $\\frac{1}{2}(1-\\frac{1}{\\nu}-o(1))$-approximation\nfor the same problem.\n\\end{theorem}\n\nIn case the matroid contains two disjoint bases ($\\nu=2$), we obtain a\n$(\\frac14-o(1))$-approximation, improving the previously known factor of $\\frac16-o(1)$ \\cite{LMNS09}.\nIn the range of $\\nu \\in (1,2]$, our positive and negative results are within a factor of $2$.\nFor maximizing a submodular function over the bases of a general matroid, \nwe obtain the following.\n\n\\begin{corollary}\n\\label{coro:matroid-bases}\nFor the problem $\\max \\{f(S): S \\in {\\cal B}\\}$, where $f(S)$ is a nonnegative submodular function,\nand ${\\cal B}$ is a collection of bases in a matroid, any constant-factor approximation requires\nan exponential number of value queries.\n\\end{corollary}\n\nWe also consider the problem of maximizing a nonnegative submodular function subject to a matroid independence\nconstraint. \n\n\\begin{theorem}\n\\label{thm:matroid-indep}\nFor any $\\epsilon>0$, a $(\\frac12 + \\epsilon)$-approximation for the problem $\\max \\{f(S): S \\in {\\cal I} \\}$,\nwhere $f(S)$ is a nonnegative submodular function, and ${\\cal I}$ is a collection of independent sets in a matroid,\nwould require exponentially many value queries.\n\nOn the other hand, there is a randomized $\\frac14 (-1+\\sqrt{5}-o(1)) \\simeq 0.309$-approximation for the same problem.\n\\end{theorem}\n\nOur algorithmic result improves a previously known $(\\frac14 - o(1))$-approximation \\cite{LMNS09}.\nThe hardness threshold follows from our general result, but also quite easily from \\cite{FMV07}.\n\n\n\\medskip\n\\noindent{\\bf Hardness from the symmetry gap.}\nNow we describe our general hardness result.\nConsider an instance $\\max \\{f(S): S \\in {\\cal F}\\}$ which exhibits a certain\ndegree of symmetry. This is formalized by the notion of a {\\em symmetry group}\n${\\cal G}$. We consider permutations $\\sigma \\in {\\bf S}(X)$ where ${\\bf S}(X)$\nis the symmetric group (of all permutations) on the ground set $X$.\nWe also use $\\sigma$ for the naturally induced mapping of subsets of $X$:\n$\\sigma(S) = \\{ \\sigma(i): i \\in S \\}$.\nWe say that the instance is invariant\nunder ${\\cal G} \\subset {\\bf S}(X) $, if for any $\\sigma \\in {\\cal G}$\nand any $S \\subseteq X$,\n$f(S) = f(\\sigma(S))$ and $S \\in {\\cal F} \\Leftrightarrow \\sigma(S) \\in {\\cal F}$.\nWe emphasize that even though we apply $\\sigma$ to sets, it must be derived\nfrom a permutation on $X$.\nFor ${\\bf x} \\in [0,1]^X$, we define the ``symmetrization of ${\\bf x}$'' as\n$$\\bar{{\\bf x}} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x})],$$\nwhere $\\sigma \\in {\\cal G}$ is uniformly random and $\\sigma({\\bf x})$ denotes ${\\bf x}$ with coordinates permuted by $\\sigma$.\n\n\\medskip\n{\\bf Erratum:}\nThe main hardness result in the conference version of this paper \\cite{Vondrak09} was formulated\nfor an arbitrary feasibility constraint ${\\cal F}$, invariant under ${\\cal G}$.\nUnfortunately, this was an error and the theorem does not hold in that form\n --- an algorithm could gather some information from querying ${\\cal F}$, and combining this\nwith information obtained by querying the objective function $f$ it could possibly determine the hidden optimal solution.\nThe possibility of gathering information from the membership oracle for ${\\cal F}$ was neglected in the proof. (The reason for this\nwas probably that the feasibility constraints used in concrete applications of the theorem were very simple and indeed\ndid not provide any information about the optimum.)\nNevertheless, to correct this issue, one needs to impose a stronger symmetry constraint on ${\\cal F}$, namely\nthe condition that $S \\in {\\cal F}$ depends only on the symmetrized version of $S$, $\\overline{{\\bf 1}_S} = {\\bf E}_{\\sigma \\in {\\cal G}}[{\\bf 1}_{\\sigma(S)}]$.\nThis is the case in all the applications of the hardness theorem in \\cite{Vondrak09} and \\cite{OV11} and hence\nthese applications are not affected.\n\n\\begin{definition}\n\\label{def:total-sym}\nWe call an instance $\\max \\{f(S): S \\in {\\cal F}\\}$ on a ground set $X$ strongly symmetric with respect to a group of permutations ${\\cal G}$ on $X$, if $f(S) = f(\\sigma(S))$ for all $S \\subseteq X$ and $\\sigma \\in {\\cal G}$, and $S \\in {\\cal F} \\Leftrightarrow S' \\in {\\cal F}$ whenever ${\\bf E}_{\\sigma \\in {\\cal G}}[{\\bf 1}_{\\sigma(S)}] = {\\bf E}_{\\sigma \\in {\\cal G}}[{\\bf 1}_{\\sigma(S')}] $.\n\\end{definition}\n\n\\noindent{\\bf Example.} A cardinality constraint, ${\\cal F} = \\{S \\subseteq [n]: |S| \\leq k \\}$, is strongly symmetric with respect to all permutations, because the condition $S \\in {\\cal F}$ depends only on the symmetrized vector $\\overline{{\\bf 1}_S} = \\frac{|S|}{n} {\\bf 1}$.\nSimilarly, a partition matroid constraint, ${\\cal F} = \\{S \\subseteq [n]: |S \\cap X_i| \\leq k \\}$ for disjoint sets $X_i$, is also strongly symmetric.\nOn the other hand, consider a family of feasible solution ${\\cal F} = \\{ \\{1,2\\}, \\{2,3\\}, \\{3,4\\}, \\{4,1\\} \\}$. This family is invariant under a group generated by the cyclic rotation $1 \\rightarrow 2 \\rightarrow 3 \\rightarrow 4 \\rightarrow 1$. It is not strongly symmetric, because the condition $S \\in {\\cal F}$ does not depend only the symmetrized vector $\\overline{{\\bf 1}_S} = (\\frac14 |S|,\\frac14 |S|,\\frac14 |S|,\\frac14 |S|)$;\nsome pairs are feasible and others are not.\n\n\n\\\n\nNext, we define the symmetry gap as the ratio between the optimal solution\nof $\\max \\{F({\\bf x}): {\\bf x} \\in P({\\cal F})\\}$ and the best {\\em symmetric} solution\nof this problem.\n\n\\begin{definition}[Symmetry gap]\nLet $\\max \\{ f(S): S \\in {\\cal F} \\}$ be an instance on a ground set $X$,\nwhich is strongly symmetric with respect to ${\\cal G} \\subset {\\bf S}(X)$.\nLet $F({\\bf x}) = {\\bf E}[f(\\hat{{\\bf x}})]$ be the multilinear extension of $f(S)$\nand $P({\\cal F}) = \\mbox{conv}(\\{{\\bf 1}_I: I \\in {\\cal F} \\})$ the polytope\nassociated with $\\cal F$. Let $\\bar{{\\bf x}} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x})]$.\nThe symmetry gap of $\\max \\{ f(S): S \\in {\\cal F} \\}$ is defined as\n$\\gamma = \\overline{OPT} \/ OPT$ where\n$$OPT = \\max \\{F({\\bf x}): {\\bf x} \\in P({\\cal F})\\},$$\n$$\\overline{OPT} = \\max \\{F(\\bar{{\\bf x}}): {\\bf x} \\in P({\\cal F}) \\}.$$\n\\end{definition}\n\n\\noindent\nWe give examples of computing the symmetry gap in Section~\\ref{section:hardness-applications}.\nNext, we need to define the notion of a {\\em refinement} of an instance.\nThis is a natural way to extend a family of feasible sets to a larger ground set.\nIn particular, this operation preserves the types of constraints that we care about,\nsuch as cardinality constraints, matroid independence, and matroid base constraints.\n\n\\begin{definition}[Refinement]\nLet ${\\cal F} \\subseteq 2^X$, $|X|=k$ and $|N| = n$. We say that\n$\\tilde{{\\cal F}}\\subseteq 2^{N \\times X}$ is a refinement of $\\cal F$, if\n$$ \\tilde{{\\cal F}} = \\left\\{ \\tilde{S} \\subseteq N \\times X \\ \\big| \\ (x_1,\\ldots,x_k) \\in P({{\\cal F}})\n \\mbox{ where } x_j = \\frac{1}{n} |\\tilde{S} \\cap (N \\times \\{j\\})| \\right\\}. $$\n\\end{definition}\n\nIn other words, in the refined instance, each element $j \\in X$ is replaced by a set $N \\times \\{j\\}$.\nWe call this set the {\\em cluster} of elements corresponding to $j$.\nA set $\\tilde{S}$ is in $\\tilde{{\\cal F}}$ if and only if the fractions $x_j$ of the respective clusters that are intersected by $\\tilde{S}$ form a vector ${\\bf x} \\in P({\\cal F})$, i.e. a convex combination of sets in ${\\cal F}$.\n\nOur main result is that the symmetry gap for any strongly symmetric instance\ntranslates automatically into hardness of approximation for refined instances.\n(See Definition~\\ref{def:total-sym} for the notion of being ``strongly symmetric\".)\nWe emphasize that this is a query-complexity lower bound,\nand hence independent of assumptions such as $P \\neq NP$.\n\n\\begin{theorem}\n\\label{thm:general-hardness}\nLet $\\max \\{ f(S): S \\in {\\cal F} \\}$ be an instance of nonnegative\n(optionally monotone) submodular maximization, strongly symmetric with respect to ${\\cal G}$,\nwith symmetry gap $\\gamma = \\overline{OPT} \/ OPT$.\nLet $\\cal C$ be the class of instances $\\max \\{\\tilde{f}(S): S \\in \\tilde{{\\cal F}}\\}$\nwhere $\\tilde{f}$ is nonnegative (optionally monotone) submodular\nand $\\tilde{{\\cal F}}$ is a refinement of ${\\cal F}$.\nThen for every $\\epsilon > 0$,\nany (even randomized) $(1+\\epsilon) \\gamma$-approximation algorithm for the class $\\cal C$ would require\nexponentially many value queries to $\\tilde{f}(S)$.\n\\end{theorem}\n\nWe remark that the result holds even if the class $\\cal C$ is restricted\nto instances which are themselves symmetric under a group related to ${\\cal G}$\n(see the discussion in Section~\\ref{section:hardness-proof}, after the proofs of Theorem~\\ref{thm:general-hardness}\nand \\ref{thm:multilinear-hardness}).\nOn the algorithmic side, submodular maximization seems easier for symmetric instances\nand in this case we obtain optimal approximation factors, up to lower-order terms\n(see Section~\\ref{section:symmetric}).\n\nOur hardness construction yields impossibility results also for solving\nthe continuous problem $\\max \\{F({\\bf x}): {\\bf x} \\in P({\\cal F}) \\}$. In the case of matroid constraints,\nthis is easy to see, because an approximation to the continuous problem gives the same\napproximation factor for the discrete problem (by pipage rounding, see Appendix~\\ref{app:pipage}).\nHowever, this phenomenon\nis more general and we can show that the value of a symmetry gap translates into\nan inapproximability result for the multilinear optimization problem under any constraint\nsatisfying a symmetry condition. \n\n\\begin{theorem}\n\\label{thm:multilinear-hardness}\nLet $\\max \\{ f(S): S \\in {\\cal F} \\}$ be an instance of nonnegative\n(optionally monotone) submodular maximization, strongly symmetric with respect to ${\\cal G}$,\nwith symmetry gap $\\gamma = \\overline{OPT} \/ OPT$.\nLet $\\cal C$ be the class of instances $\\max \\{\\tilde{f}(S): S \\in \\tilde{{\\cal F}}\\}$\nwhere $\\tilde{f}$ is nonnegative (optionally monotone) submodular\nand $\\tilde{{\\cal F}}$ is a refinement of ${\\cal F}$.\nThen for every $\\epsilon > 0$,\nany (even randomized) $(1+\\epsilon) \\gamma$-approximation algorithm for the multilinear relaxation\n$\\max \\{\\tilde{F}({\\bf x}): {\\bf x} \\in P(\\tilde{{\\cal F}})\\}$ of problems in $\\cal C$ would require\nexponentially many value queries to $\\tilde{f}(S)$.\n\\end{theorem}\n\n\n\\\n\n\\noindent{\\bf Additions to the conference version and follow-up work.}\nAn extended abstract of this work appeared in IEEE FOCS 2009 \\cite{Vondrak09}.\nAs mentioned above, the main theorem in \\cite{Vondrak09} suffers from a technical flaw.\nThis does not affect the applications, but the general theorem in \\cite{Vondrak09} is not correct.\nWe provide a corrected version of the main theorem with a complete proof\n(Theorem~\\ref{thm:general-hardness}) and we extend this hardness result to the problem of solving\n the multilinear relaxation (Theorem~\\ref{thm:multilinear-hardness}).\n\nSubsequently, further work has been done which exploits the symmetry gap concept. \nIn \\cite{OV11}, it has been proved using Theorem~\\ref{thm:general-hardness} that\nmaximizing a nonnegative submodular function subject to a matroid independence constraint\nwith a factor better than $0.478$ would require exponentially many queries. Even in the case of a cardinality constraint,\n$\\max \\{f(S): |S| \\leq k\\}$ cannot be approximated within a factor better than $0.491$ using subexponentially many\nqueries \\cite{OV11}.\nIn the case of a matroid base constraint, assuming that the fractional base packing number\nis $\\nu = \\frac{k}{k-1}$ for some $k \\geq 2$, there is no $(1-e^{-1\/k}+\\epsilon)$-approximation in the value oracle model \\cite{OV11}, improving the hardness of $(1-\\frac{1}{\\nu}+\\epsilon) = (\\frac{1}{k}+\\epsilon)$-approximation from this paper.\nThese applications are not affected by the flaw in \\cite{Vondrak09},\nand they are implied by the corrected version of Theorem~\\ref{thm:general-hardness} here.\n\nRecently \\cite{DV11}, it has been proved using the symmetry gap technique that combinatorial auctions with submodular bidders do not admit any truthful-in-expectation $1\/m^\\gamma$-approximation, where $m$ is the number of items and $\\gamma>0$ some absolute constant.\nThis is the first nontrivial hardness result for truthful-in-expectation mechanisms for combinatorial auctions;\nit separates the classes of monotone submodular functions and coverage functions,\nwhere a truthful-in-expectation $(1-1\/e)$-approximation is possible \\cite{DRY11}.\nThe proof is self-contained and does not formally refer to \\cite{Vondrak09}.\n\nMoreover, this hardness result for truthful-in-expectation mechanisms as well as the main hardness result in this paper\nhave been converted from the oracle setting to a computational complexity setting \\cite{DV12a,DV12b}. This recent work\nshows that the hardness of approximation arising from symmetry gap is not limited to instances given by an oracle, but holds\nalso for instances encoded explicitly on the input, under a suitable complexity-theoretic assumption.\n\n\\\n\n\\noindent{\\bf Organization.}\nThe rest of the paper is organized as follows.\nIn Section~\\ref{section:hardness-applications}, we present applications\nof our main hardness result (Theorem~\\ref{thm:general-hardness}) to concrete cases,\nin particular we show how it implies the hardness statements in Theorem~\\ref{thm:matroid-indep}\nand \\ref{thm:matroid-bases}.\nIn Section~\\ref{section:hardness-proof}, we present the proofs of Theorem~\\ref{thm:general-hardness}\n and Theorem~\\ref{thm:multilinear-hardness}.\nIn Section~\\ref{section:algorithms}, we prove the algorithmic results\nin Theorem~\\ref{thm:matroid-indep} and \\ref{thm:matroid-bases}.\nIn Section~\\ref{section:symmetric}, we discuss the special case of symmetric instances.\nIn the Appendix,\nwe present a few basic facts concerning submodular functions,\nan extension of pipage rounding to matroid independence polytopes (rather than matroid base polytopes),\nand other technicalities that would hinder the main exposition.\n\n\n\n\\section{From symmetry to inapproximability: applications}\n\\label{section:hardness-applications}\n\nBefore we get into the proof of Theorem~\\ref{thm:general-hardness},\nlet us show how it can be applied to a number of specific problems.\nSome of these are hardness results that were proved previously by\nan ad-hoc method. The last application is a new one\n(Theorem~\\ref{thm:matroid-bases}).\n\n\\\n\n\\noindent{\\bf Nonmonotone submodular maximization.}\nLet $X = \\{1,2\\}$ and for any $S \\subseteq X$,\n$f(S) = 1$ if $|S|=1$, and $0$ otherwise.\nConsider the instance $\\max \\{ f(S): S \\subseteq X \\}$.\nIn other words, this is the Max Cut problem on the graph $K_2$.\nThis instance exhibits a simple symmetry, the group of all (two) permutations\non $\\{1,2\\}$. We get $OPT = F(1,0) = F(0,1) = 1$, while $\\overline{OPT} = F(1\/2,1\/2)\n = 1\/2$. Hence, the symmetry gap is $1\/2$.\n\n\\begin{figure}[here]\n\\begin{tikzpicture}[scale=.50]\n\n\\draw (-10,0) node {};\n\n\\filldraw [fill=gray,line width=1mm] (0,0) rectangle (4,2);\n\n\\filldraw [fill=white] (0.5,1) .. controls +(0,1) and +(0,1) .. (1.5,1)\n .. controls +(0,-1) and +(0,-1) .. (0.5,1);\n\n\n\\fill (1,1) circle (5pt);\n\\fill (3,1) circle (5pt);\n\\draw (1,1) -- (3,1);\n\n\\draw (-1,1) node {$X$};\n\n\\draw (7,1) node {$\\bar{x}_1 = \\bar{x}_2 = \\frac{1}{2}$};\n\n\\end{tikzpicture}\n\n\\caption{Symmetric instance for nonmonotone submodular maximization: Max Cut on the graph $K_2$.\nThe white set denotes the optimal solution, while $\\bar{{\\bf x}}$ is the (unique) symmetric solution.}\n\n\\end{figure}\n\n\nSince $f(S)$ is nonnegative submodular and there is no constraint on $S \\subseteq X$,\nthis will be the case for any refinement of the instance as well.\nTheorem~\\ref{thm:general-hardness} implies immediately the following: any algorithm achieving\na $(\\frac12 + \\epsilon)$-approximation for nonnegative (nonmonotone) submodular\nmaximization requires exponentially many value queries\n(which was previously known \\cite{FMV07}). Note that a ``trivial instance\" implies\na nontrivial hardness result. This is typically the case in applications of Theorem~\\ref{thm:general-hardness}.\n\nThe same symmetry gap holds if we impose some simple constraints:\nthe problems $\\max \\{ f(S): |S| \\leq 1 \\}$ and $\\max \\{ f(S): |S| = 1 \\}$\nhave the same symmetry gap as above. Hence, the hardness threshold of $1\/2$\nalso holds for nonmonotone submodular maximization under cardinality constraints\nof the type $|S| \\leq n\/2$, or $|S| = n\/2$. This proves the hardness part of\nTheorem~\\ref{thm:matroid-indep}. This can be derived quite easily\nfrom the construction of \\cite{FMV07} as well.\n\n\\\n\n\\noindent{\\bf Monotone submodular maximization.}\nLet $X = [k]$ and $f(S) = \\min \\{|S|, 1\\}$.\nConsider the instance $\\max \\{f(S): |S| \\leq 1 \\}$.\nThis instance is invariant under all permutations on $[k]$, the symmetric group ${\\bf S}_k$.\nNote that the instance is {\\em strongly symmetric} (Def.~\\ref{def:total-sym}) with respect to ${\\bf S}_k$,\nsince the feasibility constraint $|S| \\leq 1$ depends only on the symmetrized vector\n$\\overline{{\\bf 1}_S} = (\\frac{1}{k}|S|,\n\\ldots, \\frac{1}{k}|S|)$. We get $OPT = F(1,0,\\ldots,0) = 1$, while\n$\\overline{OPT} = F(1\/k,1\/k,\\ldots,1\/k) = 1 - (1-1\/k)^k$.\n\n\\iffalse\n\\begin{figure}[here]\n\\begin{tikzpicture}[scale=.50]\n\n\\draw (-8,0) node {};\n\n\\filldraw [fill=gray,line width=1mm] (0,0) rectangle (9,2);\n\n\\filldraw [fill=white] (0.5,1) .. controls +(0,1) and +(0,1) .. (1.5,1)\n .. controls +(0,-1) and +(0,-1) .. (0.5,1);\n\n\n\\fill (1,1) circle (5pt);\n\\fill (2,1) circle (5pt);\n\\fill (3,1) circle (5pt);\n\\fill (4,1) circle (5pt);\n\\fill (5,1) circle (5pt);\n\\fill (6,1) circle (5pt);\n\\fill (7,1) circle (5pt);\n\\fill (8,1) circle (5pt);\n\n\\draw [dotted] (1,1) -- (8,1);\n\n\\draw (-1,1) node {$X$};\n\n\\draw (11,1) node {$\\bar{x}_{i} = \\frac{1}{k}$};\n\n\\end{tikzpicture}\n\n\\caption{Symmetric instance for monotone submodular maximization.}\n\n\\end{figure}\n\\fi\n\n\n\nHere, $f(S)$ is monotone submodular and any refinement of $\\cal F$ is\na set system of the type $\\tilde{{\\cal F}} = \\{S: |S| \\leq \\ell \\}$.\nBased on our theorem, this implies that any approximation better than $1 - (1-1\/k)^k$\nfor monotone submodular maximization subject\nto a cardinality constraint would require exponentially many value queries.\nSince this holds for any fixed $k$, we get the same hardness\nresult for any $\\beta > \\lim_{k \\rightarrow \\infty} (1 - (1-1\/k)^k) = 1-1\/e$\n(which was previously known \\cite{NW78}).\n\n\\iffalse\n\n\\noindent{\\bf Submodular welfare maximization.}\nLet $X = [k] \\times [k]$,\n$ {\\cal F} = \\{ S: S$ contains at most $1$ pair $(i,j)$ for each $j \\}$,\nand $f(S) = |\\{ i: \\exists (i,j) \\in S \\}|$.\nConsider the instance $\\max \\{f(S): S \\in {\\cal F}\\}$.\nThis instance can be interpreted as an allocation problem of $k$ items to\n$k$ players. A set $S$ represents an assignment in the sense that\n$(i,j) \\in S$ if item $j$ is allocated to player $i$.\nA player is satisfied is she receives at least 1 item;\nthe objective function is the number of satisfied players.\nThis instance exhibits the symmetry of all permutations of the players\n(and also all permutations of the items, although we do not use it here).\nNote that the feasibility constraint $S \\in {\\cal F}$ depends only on the symmetrized\nvector $\\overline{{\\bf 1}_S}$ which averages out the allocation of each item\nacross all players. Therefore the instance is strongly symmetric with respect\nto permutations of players.\nAn optimum solution allocates each item to a different player, and $OPT = k$.\nThe symmetrized optimum allocates a $1\/k$-fraction of each item to each player,\nwhich gives $\\overline{OPT} = F(\\frac{1}{k},\\frac{1}{k},\\ldots,\\frac{1}{k}) =\n k(1-(1-\\frac{1}{k})^k)$.\n\n\\begin{figure}[here]\n\\begin{tikzpicture}[scale=.50]\n\n\\draw (-8,0) node {};\n\n\\filldraw [fill=gray,line width=1mm] (0,0) rectangle (8,2);\n\\filldraw [fill=white] (0.5,1) .. controls +(0,1) and +(0,1) .. (1.5,1)\n .. controls +(0,-1) and +(0,-1) .. (0.5,1);\n\\filldraw [fill=gray,line width=1mm] (0,2) rectangle (8,4);\n\\filldraw [fill=white] (2.5,3) .. controls +(0,1) and +(0,1) .. (3.5,3)\n .. controls +(0,-1) and +(0,-1) .. (2.5,3);\n\\filldraw [fill=gray,line width=1mm] (0,4) rectangle (8,6);\n\\filldraw [fill=white] (4.5,5) .. controls +(0,1) and +(0,1) .. (5.5,5)\n .. controls +(0,-1) and +(0,-1) .. (4.5,5);\n\\filldraw [fill=gray,line width=1mm] (0,6) rectangle (8,8);\n\\filldraw [fill=white] (6.5,7) .. controls +(0,1) and +(0,1) .. (7.5,7)\n .. controls +(0,-1) and +(0,-1) .. (6.5,7);\n\n\n\\fill (1,1) circle (5pt);\n\\fill (3,1) circle (5pt);\n\\fill (5,1) circle (5pt);\n\\fill (7,1) circle (5pt);\n\\fill (1,3) circle (5pt);\n\\fill (3,3) circle (5pt);\n\\fill (5,3) circle (5pt);\n\\fill (7,3) circle (5pt);\n\\fill (1,5) circle (5pt);\n\\fill (3,5) circle (5pt);\n\\fill (5,5) circle (5pt);\n\\fill (7,5) circle (5pt);\n\\fill (1,7) circle (5pt);\n\\fill (3,7) circle (5pt);\n\\fill (5,7) circle (5pt);\n\\fill (7,7) circle (5pt);\n\n\\draw [dotted] (1,1) -- (7,1);\n\\draw [dotted] (1,3) -- (7,3);\n\\draw [dotted] (1,5) -- (7,5);\n\\draw [dotted] (1,7) -- (7,7);\n\n\\draw (-1.5,4) node {players};\n\\draw (4,9) node {items};\n\n\\draw (10,4) node {$\\bar{x}_{ij} = \\frac{1}{k}$};\n\n\\end{tikzpicture}\n\n\\caption{Symmetric instance for submodular welfare maximization.}\n\n\\end{figure}\n\nA refinement of this instance can be interpreted as an allocation problem\nwhere we have $n$ copies of each item, we still have $k$ players,\nand the utility functions are monotone submodular. Our theorem implies\nthat for submodular welfare maximization with $k$ players, a better approximation factor\nthan $1 - (1-\\frac{1}{k})^k$ is impossible.\\footnote{We note that \nour theorem here assumes an oracle model where only the total value can be\nqueried for a given allocation. This is actually enough for the $(1-1\/e)$-approximation\nof \\cite{Vondrak08} to work. However, the hardness result holds even if each player's valuation\nfunction can be queried separately; this result was proved in \\cite{MSV08}.}\n\n\\fi\n\n\\\n\n\\noindent{\\bf Submodular maximization over matroid bases.}\nLet $X = A \\cup B$, $A = \\{a_1, \\ldots, a_k\\}$, $B = \\{b_1, \\ldots, b_k\\}$\nand ${\\cal F} = \\{ S: |S \\cap A| = 1 \\ \\& \\ |S \\cap B| = k-1 \\}$.\nWe define $f(S) = \\sum_{i=1}^{k} f_i(S)$ where\n$f_i(S) = 1$ if $a_i \\in S \\ \\& \\ b_i \\notin S$, and $0$ otherwise.\nThis instance can be viewed as a Maximum Directed Cut problem on a graph\nof $k$ disjoint arcs, under the constraint that exactly one arc tail\nand $k-1$ arc heads should be on the left-hand side ($S$).\nAn optimal solution is for example $S = \\{a_1, b_2, b_3, \\ldots, b_k \\}$,\nwhich gives $OPT = 1$.\nThe symmetry here is that we can apply the same permutation to $A$ and $B$\nsimultaneously. Again, the feasibility of a set $S$ depends only on the symmetrized vector $\\overline{{\\bf 1}_S}$:\nin fact $S \\in {\\cal F}$ if and only if $\\overline{{\\bf 1}_S} = (\\frac1k,\\ldots,\\frac1k,1-\\frac1k,\\ldots,1-\\frac1k)$.\nThere is a unique symmetric solution\n$\\bar{{\\bf x}} = (\\frac1k,\\ldots,\\frac1k,1-\\frac1k,\\ldots,1-\\frac1k)$, and\n$\\overline{OPT} = F(\\bar{{\\bf x}}) = {\\bf E}[f(\\hat{\\bar{{\\bf x}}})] = \\sum_{i=1}^{k}\n {\\bf E}[f_i(\\hat{\\bar{{\\bf x}}})] = \\frac1k$\n(since each arc appears in the directed cut induced by $\\hat{\\bar{{\\bf x}}}$ with probability $\\frac{1}{k^2}$).\n\n\\begin{figure}[here]\n\\begin{tikzpicture}[scale=.60]\n\n\\draw (-6,0) node {};\n\n\\filldraw [fill=gray,line width=1mm] (0,0) rectangle (9,2);\n\\filldraw [fill=gray,line width=1mm] (0,2) rectangle (9,4);\n\n\\filldraw [fill=white] (0.5,3) .. controls +(0,1) and +(0,1) .. (1.5,3)\n .. controls +(0,-1) and +(0,-1) .. (0.5,3);\n\\filldraw [fill=white] (1.5,1) .. controls +(0,1) and +(0,1) .. (8.5,1)\n .. controls +(0,-1) and +(0,-1) .. (1.5,1);\n\n\n\\fill (1,1) circle (5pt);\n\\fill (2,1) circle (5pt);\n\\fill (3,1) circle (5pt);\n\\fill (4,1) circle (5pt);\n\\fill (5,1) circle (5pt);\n\\fill (6,1) circle (5pt);\n\\fill (7,1) circle (5pt);\n\\fill (8,1) circle (5pt);\n\\fill (1,3) circle (5pt);\n\\fill (2,3) circle (5pt);\n\\fill (3,3) circle (5pt);\n\\fill (4,3) circle (5pt);\n\\fill (5,3) circle (5pt);\n\\fill (6,3) circle (5pt);\n\\fill (7,3) circle (5pt);\n\\fill (8,3) circle (5pt);\n\\draw[-latex] [line width=0.5mm] (1,3) -- (1,1);\n\\draw[-latex] [line width=0.5mm] (2,3) -- (2,1);\n\\draw[-latex] [line width=0.5mm] (3,3) -- (3,1);\n\\draw[-latex] [line width=0.5mm] (4,3) -- (4,1);\n\\draw[-latex] [line width=0.5mm] (5,3) -- (5,1);\n\\draw[-latex] [line width=0.5mm] (6,3) -- (6,1);\n\\draw[-latex] [line width=0.5mm] (7,3) -- (7,1);\n\\draw[-latex] [line width=0.5mm] (8,3) -- (8,1);\n\n\n\\draw (-1,1) node {$B$};\n\\draw (-1,3) node {$A$};\n\n\\draw (11,3) node {$\\bar{x}_{a_i} = \\frac{1}{k}$};\n\\draw (11,1) node {$\\bar{x}_{b_i} = 1 - \\frac{1}{k}$};\n\n\\end{tikzpicture}\n\n\\caption{Symmetric instance for submodular maximization over matroid bases.}\n\n\\end{figure}\nThe refined instances are instances of (nonmonotone) submodular maximization\nover the bases of a matroid, where the ground set is partitioned into $A \\cup B$\nand we should take a $\\frac1k$-fraction of $A$ and a $(1-\\frac1k)$-fraction of $B$.\n(This means that the fractional packing number of bases is $\\nu = \\frac{k}{k-1}$.)\nOur theorem implies that for this class of instances, an approximation better\nthan $1\/k$ is impossible - this proves the hardness part of Theorem~\\ref{thm:matroid-bases}.\n\n\\\n\nObserve that in all the cases mentioned above, the multilinear relaxation is equivalent\nto the original problem, in the sense that any fractional solution can be rounded\nwithout any loss in the objective value. This implies that the same hardness factors apply to\nsolving the multilinear relaxation of the respective problems. In particular, using the last result\n(for matroid bases), we obtain that the multilinear optimization problem $\\max \\{ F({\\bf x}): {\\bf x} \\in P \\}$\ndoes not admit a constant factor for nonnegative submodular functions and matroid base polytopes.\n(We remark that a $(1-1\/e)$-approximation can be achieved for any {\\em monotone} submodular function\nand any solvable polytope, i.e.~polytope over which we can optimize linear functions \\cite{Vondrak08}.)\n\nAs Theorem~\\ref{thm:multilinear-hardness} shows, this holds more generally - any symmetry gap\nconstruction gives an inapproximability result for solving the multilinear optimization problem\n$\\max \\{F({\\bf x}): {\\bf x} \\in P\\}$. This in fact implies limits on what hardness results we can\npossibly hope for using this technique. For instance, we cannot prove using the symmetry gap\nthat the monotone submodular maximization problem subject to the intersection of $k$ matroid constraints\ndoes not admit a constant factor - because we would also prove that the respective multilinear relaxation\ndoes not admit such an approximation. But we know from \\cite{Vondrak08} that a $(1-1\/e)$-approximation\nis possible for the multilinear problem in this case.\n\nHence, the hardness arising from the symmetry gap is related to the difficulty of solving\nthe multilinear optimization problem rather than the difficulty of rounding a fractional solution.\nThus this technique is primarily suited to optimization problems where the multilinear optimization problem\ncaptures closely the original discrete problem.\n \n\n\n\n\\section{From symmetry to inapproximability: proof}\n\\label{section:hardness-proof}\n\n\\paragraph{The roadmap}\nAt a high level,\nour proof resembles the constructions of \\cite{FMV07,MSV08}.\nWe construct instances based on continuous functions $F({\\bf x})$, $G({\\bf x})$,\nwhose optima differ by a gap\nfor which we want to prove hardness. Then we show that after a certain perturbation,\nthe two instances are very hard to distinguish.\nThis paper generalizes the ideas of \\cite{FMV07,MSV08} and brings\ntwo new ingredients. First, we show that the functions\n$F({\\bf x}), G({\\bf x})$, which are ``pulled out of the hat'' in \\cite{FMV07,MSV08},\ncan be produced in a natural way from the multilinear relaxation\nof the respective problem, using the notion of a {\\em symmetry gap}.\nSecondly, the functions $F({\\bf x}), G({\\bf x})$ are perturbed in a way that makes\nthem indistinguishable and this forms the main technical part\nof the proof. In \\cite{FMV07}, this step is quite simple.\nIn \\cite{MSV08}, the perturbation is more complicated, but still relies\non properties of the functions $F({\\bf x}), G({\\bf x})$ specific to that application.\nThe construction that we present here (Lemma~\\ref{lemma:final-fix})\nuses the symmetry properties of a fixed instance in a generic fashion.\n\n\n\\\n\nFirst, let us present an outline of our construction. Given an instance\n$\\max \\{f(S): S \\in {\\cal F}\\}$ exhibiting a symmetry gap $\\gamma$,\nwe consider two smooth submodular\\footnote{\"Smooth submodularity\"\nmeans the condition $\\mixdiff{F}{x_i}{x_j} \\leq 0$ for all $i,j$.}\nfunctions, $F({\\bf x})$ and $G({\\bf x})$.\nThe first one is the multilinear extension $F({\\bf x}) = {\\bf E}[f(\\hat{{\\bf x}})]$,\nwhile the second one is its symmetrized version $G({\\bf x}) = F(\\bar{{\\bf x}})$.\nWe modify these functions slightly so that we obtain\nfunctions $\\hat{F}({\\bf x})$ and $\\hat{G}({\\bf x})$ with the following property:\nFor any vector ${\\bf x}$ which is close to its symmetrized version\n$\\bar{{\\bf x}}$, $\\hat{F}({\\bf x}) = \\hat{G}({\\bf x})$.\nThe functions $\\hat{F}({\\bf x}), \\hat{G}({\\bf x})$ induce instances of\nsubmodular maximization on the refined ground sets. The way we define\ndiscrete instances based on $\\hat{F}({\\bf x}), \\hat{G}({\\bf x})$ is natural,\nusing the following lemma.\nEssentially, we interpret the fractional variables as fractions of clusters\nin the refined instance.\n\n\\begin{lemma}\n\\label{lemma:smooth-submodular}\nLet $F:[0,1]^X \\rightarrow {\\boldmath R}$,\n $N = [n]$, $n \\geq 1$, and define $f:2^{N \\times X} \\rightarrow {\\boldmath R}$\nso that $f(S) = F({\\bf x})$ where $x_i = \\frac{1}{n} |S \\cap (N \\times \\{i\\})|$. Then\n\\begin{enumerate}\n\\item If $\\partdiff{F}{x_i} \\geq 0$ everywhere for each $i$, then\n$f$ is monotone.\n\\item If the first partial derivatives of $F$ are absolutely continuous\\footnote{\nA function $F:[0,1]^X \\rightarrow {\\boldmath R}$ is absolutely continuous,\nif $\\forall \\epsilon>0; \\exists \\delta>0; \\sum_{i=1}^{t} ||{\\bf x}_i-{\\bf y}_i|| < \\delta \\Rightarrow\n\\sum_{i=1}^{t} |F({\\bf x}_i) - F({\\bf y}_i)| < \\epsilon$.} and $\\mixdiff{F}{x_i}{x_j} \\leq 0$ almost everywhere for all $i,j$,\nthen $f$ is submodular.\n\\end{enumerate}\n\\end{lemma}\n\n\\begin{proof}\nFirst, assume $\\partdiff{F}{x_i} \\geq 0$ everywhere for all $i$. This implies that $F$ is nondecreasing in every coordinate,\ni.e. $F({\\bf x}) \\leq F({\\bf y})$ whenever ${\\bf x} \\leq {\\bf y}$. This means that $f(S) \\leq f(T)$ whenever $S \\subseteq T$.\n\nNext, assume $\\partdiff{F}{x_i}$ is absolutely continuous for each $i$ and $\\mixdiff{F}{x_j}{x_i} \\leq 0$ almost everywhere for all $i,j$.\nWe want to prove that $\\partdiff{F}{x_i} |_{\\bf x} \\geq \\partdiff{F}{x_i} |_{\\bf y}$ whenever ${\\bf x} \\leq {\\bf y}$, which implies that the marginal\nvalues of $f$ are nonincreasing.\n\nLet ${\\bf x} \\leq {\\bf y}$,\nfix $\\delta>0$ arbitrarily small, and pick ${\\bf x}',{\\bf y}'$ such that $||{\\bf x}'-{\\bf x}||<\\delta, ||{\\bf y}'-{\\bf y}||<\\delta, {\\bf x}' \\leq {\\bf y}'$ and\non the line segment $[{\\bf x}', {\\bf y}']$, we have $\\mixdiff{F}{x_j}{x_i} \\leq 0$ except for a set of (1-dimensional) measure zero. If such a pair of points\n${\\bf x}', {\\bf y}'$ does not exist, it means that there are sets $A,B$ of positive measure such that\n${\\bf x} \\in A, {\\bf y} \\in B$ and for any\n${\\bf x}' \\in A, {\\bf y}' \\in B$, the line segment $[{\\bf x}',{\\bf y}']$ contains a subset of positive (1-dimensional) measure where $\\mixdiff{F}{x_j}{x_i}$ \nfor some $j$ is positive or undefined. This would imply that $[0,1]^X$ contains a subset of positive measure where\n$\\mixdiff{F}{x_j}{x_i}$ for some $j$ is positive or undefined, which we assume is not the case.\n\nTherefore, there is a pair of points ${\\bf x}', {\\bf y}'$ as described above.\nWe compare $\\partdiff{F}{x_i} |_{{\\bf x}'}$ and $\\partdiff{F}{x_i} |_{{\\bf y}'}$\nby integrating along the line segment $[{\\bf x}', {\\bf y}']$. Since $\\partdiff{F}{x_i}$ is absolutely continuous and $\\partdiff{}{x_j} \\partdiff{F}{x_i} = \\mixdiff{F}{x_j}{x_i} \\leq 0$ along this line segment for all $j$ except for a set of measure zero, we obtain $\\partdiff{F}{x_i} |_{{\\bf x}'} \\geq \\partdiff{F}{x_i} |_{{\\bf y}'}$. This is true for ${\\bf x}', {\\bf y}'$ arbitrarily\nclose to ${\\bf x}, {\\bf y}$, and hence by continuity of $\\partdiff{F}{x_i}$, we get $\\partdiff{F}{x_i} |_{{\\bf x}} \\geq \\partdiff{F}{x_i} |_{{\\bf y}}$.\nThis implies that the marginal values of $f$ are nonincreasing.\n\\end{proof}\n\n\nThe way we construct $\\hat{F}({\\bf x}), \\hat{G}({\\bf x})$ is such that,\ngiven a large enough refinement of the ground set,\nit is impossible to distinguish the instances corresponding to\n$\\hat{F}({\\bf x})$ and $\\hat{G}({\\bf x})$. As we argue more precisely later,\nthis holds because if the ground set is large and labeled in a random way\n(considering the symmetry group of the instance), a query about a vector ${\\bf x}$\neffectively becomes a query about the symmetrized vector $\\bar{{\\bf x}}$.\nWe would like this property to imply that all queries with high probability fall\nin the region where $\\hat{F}({\\bf x}) = \\hat{G}({\\bf x})$ and the inability to distinguish between $\\hat{F}$\nand $\\hat{G}$ gives the hardness result that we want.\nThe following lemma gives the precise properties of $\\hat{F}({\\bf x})$\nand $\\hat{G}({\\bf x})$ that we need.\n\n\\begin{lemma}\n\\label{lemma:final-fix}\nConsider a function $f:2^X \\rightarrow {\\boldmath R}_+$ invariant under a group\nof permutations $\\cal G$ on the ground set $X$.\nLet $F({\\bf x}) = {\\bf E}[f(\\hat{{\\bf x}})]$, $\\bar{{\\bf x}} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x})]$,\nand fix any $\\epsilon > 0$.\nThen there is $\\delta > 0$ and functions $\\hat{F}, \\hat{G}:[0,1]^X \\rightarrow {\\boldmath R}_+$\n(which are also symmetric with respect to ${\\cal G}$), satisfying:\n\\begin{enumerate}\n\\item For all ${\\bf x} \\in [0,1]^X$, $\\hat{G}({\\bf x}) = \\hat{F}(\\bar{{\\bf x}})$.\n\\item For all ${\\bf x} \\in [0,1]^X$, $|\\hat{F}({\\bf x}) - F({\\bf x})| \\leq \\epsilon$.\n\\item Whenever $||{\\bf x} - \\bar{{\\bf x}}||^2 \\leq \\delta$, $\\hat{F}({\\bf x}) = \\hat{G}({\\bf x})$\nand the value depends only on $\\bar{{\\bf x}}$.\n\\item The first partial derivatives of $\\hat{F}, \\hat{G}$ are absolutely continuous. \n\\item If $f$ is monotone, then $\\partdiff{\\hat{F}}{x_i} \\geq 0$ and\n $\\partdiff{\\hat{G}}{x_i} \\geq 0$ everywhere.\n\\item If $f$ is submodular, then $\\mixdiff{\\hat{F}}{x_i}{x_j} \\leq 0$ and\n $\\mixdiff{\\hat{G}}{x_i}{x_j} \\leq 0$ almost everywhere.\n\\end{enumerate}\n\\end{lemma}\n\nThe proof of this lemma is the main technical part of this paper\nand we defer it to the end of this section. Assuming this lemma,\nwe first finish the proof of the main theorem. We prove the following.\n\n\\begin{lemma}\n\\label{lemma:indistinguish}\nLet $\\hat{F}, \\hat{G}$ be the two functions provided by Lemma~\\ref{lemma:final-fix}.\nFor a parameter $n \\in {\\boldmath Z}_+$ and $N = [n]$, define two discrete functions\n$\\hat{f}, \\hat{g}: 2^{N \\times X} \\rightarrow {\\boldmath R}_+$ as follows:\nLet $\\sigma^{(i)}$ be an arbitrary permutation in ${\\cal G}$ for each $i \\in N$.\nFor every set $S \\subseteq N \\times X$, we define a vector $\\xi(S) \\in [0,1]^X$ by\n$$ \\xi_j(S) = \\frac{1}{n} \\left|\\{i \\in N: (i,\\sigma^{(i)}(j)) \\in S \\}\\right|.$$\nLet us define:\n$$ \\hat{f}(S) = \\hat{F}(\\xi(S)), \\ \\ \\ \\ \\ \\hat{g}(S) = \\hat{G}(\\xi(S)).$$\nIn addition, let $\\tilde{{\\cal F}} = \\{\\tilde{S}: \\xi(\\tilde{S}) \\in P({\\cal F})\\}$\nbe a feasibility constraint such that the condition $S \\in {\\cal F}$ depends\nonly on the symmetrized vector $\\overline{{\\bf 1}_S}$.\nThen deciding whether an instance given by value\/membership oracles is\n$\\max \\{\\hat{f}(S): S \\in \\tilde{{\\cal F}}\\}$ or $\\max \\{\\hat{g}(S): S \\in \\tilde{{\\cal F}} \\}$\n(even by a randomized algorithm, with any constant probability of success)\nrequires an exponential number of queries.\n\\end{lemma}\n\n\\begin{proof}\nLet $\\sigma^{(i)} \\in {\\cal G}$ be chosen independently at random for each $i \\in N$ and consider the instances\n$\\max\\{\\hat{f}(S): S \\in \\tilde{{\\cal F}} \\}$, $\\max \\{\\hat{g}(S): S \\in \\tilde{{\\cal F}} \\}$ as described in the lemma. We show that\nevery deterministic algorithm will follow the same computation path and return the same answer on both instances,\nwith high probability. By Yao's principle, this means that every randomized algorithm returns\nthe same answer for the two instances with high probability, for some particular $\\sigma^{(i)} \\in {\\cal G}$.\n\nThe feasible sets in the refined instance, $S \\in \\tilde{{\\cal F}}$, are such that\nthe respective vector $\\xi(S)$ is in the polytope $P({{\\cal F}})$.\nSince the instance is strongly symmetric, the condition $S \\in {\\cal F}$ depends\nonly on the symmetrized vector $\\overline{{\\bf 1}_S}$. Hence, the condition $\\xi(S) \\in P({\\cal F})$\ndepends only on the symmetrized vector $\\overline{\\xi(S)}$. Therefore, $S \\in \\tilde{{\\cal F}} \\Leftrightarrow\n\\xi(S) \\in P({\\cal F}) \\Leftrightarrow \\overline{\\xi(S)} \\in P({\\cal F})$.\nWe have $\\overline{\\xi(S)}_j = {\\bf E}_{\\sigma \\in {\\cal G}}[\\frac{1}{n} \\left|\\{i \\in N: (i,\\sigma^{(i)}(\\sigma(j))) \\in S \\}\\right|].$ The distribution of $\\sigma^{(i)} \\circ \\sigma$ is the same as that of $\\sigma$, i.e.~uniform over ${\\cal G}$. Hence, $\\overline{\\xi(S)}_j$ and consequently the condition $S \\in \\tilde{{\\cal F}}$ does not depend on $\\sigma^{(i)}$ in any way.\nIntuitively, an algorithm cannot learn any information about the permutations $\\sigma^{(i)}$ by querying the feasibility oracle,\nsince the feasibility condition $S \\in \\tilde{{\\cal F}}$ does not depend on $\\sigma^{(i)}$ for any $i \\in N$.\n\nThe main part of the proof is to show that even queries to the objective function are unlikely to reveal\nany information about the permutations $\\sigma^{(i)}$.\nThe key observation is that for any fixed query $Q$ to the objective function,\nthe associated vector ${\\bf q} = \\xi(Q)$ is very likely to be close\nto its symmetrized version $\\bar{{\\bf q}}$.\nTo see this, consider a query $Q$. The associated vector ${\\bf q} = \\xi(Q)$\nis determined by\n$$ q_j = \\frac{1}{n} |\\{i \\in N: (i,\\sigma^{(i)}(j)) \\in Q\\}|\n = \\frac{1}{n} \\sum_{i=1}^{n} Q_{ij} $$\nwhere $Q_{ij}$ is an indicator variable of the event $(i,\\sigma^{(i)}(j)) \\in Q$.\nThis is a random event due to the randomness in $\\sigma^{(i)}$.\nWe have\n$$ {\\bf E}[Q_{ij}] = \\Pr[Q_{ij}=1] = \\Pr_{\\sigma^{(i)} \\in {\\cal G}}[(i,\\sigma^{(i)}(j)) \\in Q].$$\nAdding up these expectations over $i \\in N$, we get\n$$ \\sum_{i \\in N} {\\bf E}[Q_{ij}]\n = \\sum_{i \\in N} \\Pr_{\\sigma^{(i)} \\in {\\cal G}}[(i,\\sigma^{(i)}(j)) \\in Q] \n = {\\bf E}_{\\sigma \\in {\\cal G}} [|\\{i \\in N: (i,\\sigma(j)) \\in Q \\}|].$$\nFor the purposes of expectation, the independence of $\\sigma^{(1)},\n \\ldots, \\sigma^{(n)}$ is irrelevant and that is why we replace them by one random permutation $\\sigma$.\nOn the other hand, consider the symmetrized vector $\\bar{{\\bf q}}$, with coordinates\n$$ \\bar{q}_j = {\\bf E}_{\\sigma \\in {\\cal G}}[q_{\\sigma(j)}]\n = \\frac{1}{n} {\\bf E}_{\\sigma \\in {\\cal G}}[|\\{i \\in N: (i,\\sigma^{(i)}(\\sigma(j))) \\in Q\\}|]\n = \\frac{1}{n} {\\bf E}_{\\sigma \\in {\\cal G}}[|\\{i \\in N: (i,\\sigma(j)) \\in Q \\}|] $$\nusing the fact that the distribution of $\\sigma^{(i)} \\circ \\sigma$\nis the same as the distribution of $\\sigma$ - uniformly random over ${\\cal G}$.\nNote that the vector ${\\bf q}$ depends on the random permutations $\\sigma^{(i)}$\nbut the symmetrized vector $\\bar{{\\bf q}}$ does not; this will be also useful\nin the following. For now, we summarize that\n$$ \\bar{q}_j = \\frac{1}{n} \\sum_{i=1}^{n} {\\bf E}[Q_{ij}] = {\\bf E}[q_j].$$\nSince each permutation $\\sigma^{(i)}$ is chosen independently, the random variables\n$\\{Q_{ij}: 1 \\leq i \\leq n\\}$ are independent (for a fixed $j$).\nWe can apply Chernoff's bound (see e.g. \\cite{AlonSpencer}, Corollary A.1.7):\n$$ \\Pr\\left[\\left|\\sum_{i=1}^{n} Q_{ij}- \\sum_{i=1}^{n} {\\bf E}[Q_{ij}] \\right| > a \\right]\n < 2e^{-2a^2 \/ n}.$$\nUsing $q_j = \\frac{1}{n} \\sum_{i=1}^{n} Q_{ij}$,\n $\\bar{q}_j = \\frac{1}{n} \\sum_{i=1}^{n} {\\bf E}[Q_{ij}]$\nand setting $a = n \\sqrt{\\delta\/|X|}$, we obtain\n$$ \\Pr\\left[|q_j - \\bar{q}_j| > \\sqrt{{\\delta}\/{|X|}}\\right]\n < 2e^{-2 n \\delta \/ |X|}.$$\nBy the union bound,\n\\begin{equation}\n\\label{eq:D(q)}\n \\Pr[||{\\bf q}-\\bar{{\\bf q}}||^2 > \\delta] \\leq \\sum_{j \\in X} \\Pr[|q_j-\\bar{q}_j|^2 > {\\delta}\/{|X|}] < 2|X| e^{-2n \\delta\/|X|}.\n\\end{equation}\nNote that while $\\delta$ and $|X|$ are constants, $n$ grows as the size\nof the refinement and hence the probability is exponentially small in the\nsize of the ground set $N \\times X$.\n\nDefine $D({\\bf q}) = ||{\\bf q}-\\bar{{\\bf q}}||^2$. As long as $D({\\bf q}) \\leq \\delta$ for every query issued by the algorithm,\nthe answers do not depend on the randomness of the input. This is because\nthen the values of $\\hat{F}({\\bf q})$ and $\\hat{G}({\\bf q})$ depend only on $\\bar{{\\bf q}}$,\nwhich is independent of the random permutations $\\sigma^{(i)}$,\nas we argued above.\nTherefore, assuming that $D({\\bf q}) \\leq \\delta$ for each query,\nthe algorithm will always follow the same path\nof computation and issue the same sequence of queries $\\cal S$.\n(Note that this is just a fixed sequence which can be written down\nbefore we started running the algorithm.)\nAssume that $|{\\cal S}| = e^{o(n)}$, i.e.,~the number of queries is subexponential in $n$.\nBy (\\ref{eq:D(q)}), using a union bound over all $Q \\in {\\cal S}$,\nit happens with probability $1 - e^{-\\Omega(n)}$ that ${D}({\\bf q}) = ||{\\bf q}-\\bar{{\\bf q}}||^2 \\leq \\delta$ for all $Q \\in {\\cal S}$.\n(Note that $\\delta$ and $|X|$ are constants here.)\nThen, the algorithm indeed follows this path of computation and gives the same answer.\nIn particular, the answer does not depend on whether the instance is $\\max \\{\\hat{f}(S): S \\in \\tilde{{\\cal F}}\\}$\nor $\\max \\{\\hat{g}(S): S \\in \\tilde{{\\cal F}}\\}$.\n\\end{proof}\n\n\\begin{proof}[Proof of Theorem~\\ref{thm:general-hardness}]\nFix an $\\epsilon > 0$.\nGiven an instance $\\max \\{f(S): S \\in {\\cal F}\\}$\nstrongly symmetric under ${\\cal G}$, let $\\hat{F}, \\hat{G}: [0,1]^X \\rightarrow {\\boldmath R}$\nbe the two functions provided by Lemma~\\ref{lemma:final-fix}.\nWe choose a large number $n$ and consider a refinement $\\tilde{{\\cal F}}$\non the ground set $N \\times X$, where $N = [n]$.\nWe define discrete instances of submodular maximization\n$\\max \\{ \\hat{f}(S): S \\in \\tilde{{\\cal F}} \\}$ and $\\max \\{ \\hat{g}(S): S \\in \\tilde{{\\cal F}} \\}$.\nAs in Lemma~\\ref{lemma:indistinguish}, for each $i \\in N$\nwe choose a random permutation $\\sigma^{(i)} \\in {\\cal G}$.\nThis can be viewed as a random shuffle of the labeling of the ground set\nbefore we present it to an algorithm.\nFor every set $S \\subseteq N \\times X$, we define a vector $\\xi(S) \\in [0,1]^X$ by\n$$ \\xi_j(S) = \\frac{1}{n} \\left|\\{i \\in N: (i,\\sigma^{(i)}(j)) \\in S \\}\\right|.$$\nIn other words, $\\xi_j(S)$ measures the fraction of copies of element $j$\ncontained in $S$; however, for each $i$ the $i$-copies of all elements\nare shuffled by $\\sigma^{(i)}$. Next, we define\n$$ \\hat{f}(S) = \\hat{F}(\\xi(S)), \\ \\ \\ \\ \\ \\hat{g}(S) = \\hat{G}(\\xi(S)).$$\nWe claim that $\\hat{f}$ and $\\hat{g}$ are submodular (for any fixed $\\xi$).\nNote that the effect\nof $\\sigma^{(i)}$ is just a renaming (or shuffling) of the elements\nof $N \\times X$, and hence for the purpose of proving submodularity we can assume\nthat $\\sigma^{(i)}$ is the identity for all $i$. Then, $\\xi_j(S) = \\frac{1}{n}\n |S \\cap (N \\times \\{j\\})|$. Due to Lemma~\\ref{lemma:smooth-submodular},\nthe property $\\mixdiff{\\hat{F}}{x_i}{x_j} \\leq 0$ (almost everywhere) implies that $\\hat{f}$ is submodular.\nIn addition, if the original instance was monotone, then $\\partdiff{\\hat{F}}{x_j} \\geq 0$\nand $\\hat{f}$ is monotone. The same holds for $\\hat{g}$.\n\nThe value of $\\hat{g}(S)$ for any feasible solution $S \\in \\tilde{{\\cal F}}$\nis bounded by $\\hat{g}(S) = \\hat{G}(\\xi(S)) = \\hat{F}(\\overline{\\xi(S)})\n\\leq \\overline{OPT} + \\epsilon$.\nOn the other hand, let ${\\bf x}^*$ denote a point where the optimum of the continuous problem\n$\\max \\{\\hat{F}({\\bf x}): {\\bf x} \\in P({{\\cal F}}) \\}$ is attained, i.e. $\\hat{F}({\\bf x}^*) \\geq OPT - \\epsilon$.\nFor a large enough $n$, we can approximate the point ${\\bf x}^*$ arbitrarily closely\nby a rational vector with $n$ in the denominator,\nwhich corresponds to a discrete solution $S^* \\in \\tilde{{\\cal F}}$ whose\nvalue $\\hat{f}(S^*)$ is at least, say, $OPT - 2 \\epsilon$.\nHence, the ratio between the optima of the {\\em discrete} optimization\nproblems $\\max \\{\\hat{f}(S): S \\in \\tilde{{\\cal F}}\\}$ and $\\max \\{\\hat{g}(S): S \\in \\tilde{{\\cal F}} \\}$\ncan be made at most $\\frac{\\overline{OPT} + \\epsilon}{OPT - 2 \\epsilon}$, i.e. arbitrarily close\nto the symmetry gap $\\gamma = \\frac{\\overline{OPT}}{OPT}$.\n\nBy Lemma~\\ref{lemma:indistinguish}, distinguishing the two instances\n$\\max \\{\\hat{f}(S): S \\in \\tilde{{\\cal F}} \\}$ and $\\max \\{\\hat{g}(S): S \\in \\tilde{{\\cal F}} \\}$,\neven by a randomized algorithm, requires an exponential number of value queries.\nTherefore, we cannot estimate the optimum within a factor better than $\\gamma$.\n\\end{proof}\n\nNext, we prove Theorem~\\ref{thm:multilinear-hardness} (again assuming Lemma~\\ref{lemma:final-fix}),\ni.e.~an analogous hardness result for solving the multilinear optimization problem.\n\n\n\\begin{proof}[Proof of Theorem~\\ref{thm:multilinear-hardness}]\nGiven a symmetric instance $\\max \\{f(S): S \\in {\\cal F}\\}$ and $\\epsilon>0$, we construct\nrefined and modified instances $\\max \\{\\hat{f}(S): S \\in \\tilde{{\\cal F}} \\}$,\n $\\max \\{\\hat{g}(S): S \\in \\tilde{{\\cal F}} \\}$,\nderived from the continuous functions $\\hat{F}, \\hat{G}$ provided by Lemma~\\ref{lemma:final-fix},\nexactly as we did in the proof of Theorem~\\ref{thm:general-hardness}.\nLemma~\\ref{lemma:indistinguish} states that these two instances cannot be distinguished using a subexponential number of value queries. Furthermore, the gap between the two modified instances corresponds to the symmetry gap $\\gamma$ of the original instance: $\\max \\{\\hat{f}(S): S \\in \\tilde{{\\cal F}} \\} \\geq OPT - 2 \\epsilon$ and $\\max \\{\\hat{g}(S): S \\in \\tilde{{\\cal F}} \\} \\leq \\overline{OPT} + \\epsilon = \\gamma OPT + \\epsilon$.\n\nNow we consider the multilinear relaxations of the two refined instances, $\\max \\{\\check{F}({\\bf x}): {\\bf x} \\in P(\\tilde{{\\cal F}})\\}$\nand $\\max \\{\\check{G}({\\bf x}): {\\bf x} \\in P(\\tilde{{\\cal F}})\\}$. Note that $\\check{F}, \\check{G}$, (although related to $\\hat{F},\n\\hat{G}$) are not exactly the same as the functions $\\hat{F}, \\hat{G}$; in particular, they are defined\non a larger (refined) domain. However, we show that the gap between the optima of the two instances remains\nthe same.\n\nFirst, the value of $\\max \\{\\check{F}({\\bf x}): {\\bf x} \\in P({\\cal F})\\}$ is at least the optimum of the discrete problem,\n$\\max \\{\\hat{f}(S): S \\in \\tilde{{\\cal F}} \\}$, which is at least $OPT - 2\\epsilon$ as in the proof of\nTheorem~\\ref{thm:general-hardness}. The value of $\\max \\{\\check{G}({\\bf x}): {\\bf x} \\in P({\\cal F})\\}$ can be analyzed as follows. For any fractional solution ${\\bf x} \\in P({\\cal F})$, the value of $\\check{G}({\\bf x})$\nis the expectation ${\\bf E}[\\hat{g}(\\hat{{\\bf x}})]$, where $\\hat{{\\bf x}}$ is obtained by independently rounding the\ncoordinates of ${\\bf x}$ to $\\{0,1\\}$. Recall that $\\hat{g}$ is obtained by discretizing the continuous function\n$\\hat{G}$ (using Lemma~\\ref{lemma:smooth-submodular}). In particular, $\\hat{g}(S) = \\hat{G}(\\tilde{{\\bf x}})$ where\n$\\tilde{x}_i = \\frac{1}{n} |S \\cap (N \\times \\{i\\})|$ is the fraction of the respective cluster contained in $S$, and $|N| = n$ is the size of each cluster (the refinement parameter). If ${\\bf 1}_S = \\hat{{\\bf x}}$, i.e. $S$ is chosen by independent sampling with probabilities according to ${\\bf x}$, then for large $n$ the fractions $\\frac{1}{n}|S \\cap (N \\times \\{i\\})|$ will be strongly concentrated around their expectation. As $\\hat{G}$ is continuous, we get\n$\\lim_{n \\rightarrow \\infty} {\\bf E}[\\hat{g}(\\hat{{\\bf x}})] = \\lim_{n \\rightarrow \\infty} {\\bf E}[\\hat{G}(\\tilde{{\\bf x}})]\n = \\hat{G}({\\bf E}[\\tilde{{\\bf x}}]) = \\hat{G}(\\bar{{\\bf x}})$.\nHere, $\\bar{{\\bf x}}$ is the vector ${\\bf x}$ projected back to the original ground set $X$, where the coordinates\nof each cluster have been averaged. By construction of the refinement, if ${\\bf x} \\in P(\\tilde{{\\cal F}})$ then\n$\\bar{{\\bf x}}$ is in the polytope corresponding to the original instance, $P({\\cal F})$.\nTherefore, $\\hat{G}(\\bar{{\\bf x}}) \\leq\n\\max \\{\\hat{G}({\\bf x}): {\\bf x} \\in P({\\cal F}) \\} \\leq \\gamma OPT + \\epsilon$. For large enough $n$, this means that\n$\\max \\{\\check{G}({\\bf x}): {\\bf x} \\in P(\\tilde{{\\cal F}}) \\} \\leq \\gamma OPT + 2 \\epsilon$. This holds for\nan arbitrarily small fixed $\\epsilon>0$, and hence the gap between the instances \n$\\max \\{\\check{F}({\\bf x}): {\\bf x} \\in P(\\tilde{{\\cal F}})\\}$ and $\\max \\{\\check{G}({\\bf x}): {\\bf x} \\in P(\\tilde{{\\cal F}})\\}$\n(which cannot be distinguished) can be made arbitrarily close to $\\gamma$.\n\\end{proof}\n\n\n\\paragraph{Hardness for symmetric instances}\nWe remark that since Lemma~\\ref{lemma:final-fix} provides functions $\\hat{F}$ and $\\hat{G}$ symmetric\nunder ${\\cal G}$, the refined instances that we define are invariant with respect to\nthe following symmetries: permute the copies of each element in an arbitrary\nway, and permute the classes of copies according to any permutation\n$\\sigma \\in {\\cal G}$. This means that our hardness results also hold\nfor instances satisfying such symmetry properties.\n\n\\\n\nIt remains to prove Lemma~\\ref{lemma:final-fix}.\nBefore we move to the final construction of $\\hat{F}({\\bf x})$ and $\\hat{G}({\\bf x})$,\nwe construct as an intermediate step a function $\\tilde{F}({\\bf x})$ which is helpful\nin the analysis.\n\n\\\n\n\\paragraph{Construction}\nLet us construct a function $\\tilde{F}({\\bf x})$ which satisfies the following:\n\\begin{itemize}\n\\item For ${\\bf x}$ ``sufficiently close'' to $\\bar{{\\bf x}}$, $\\tilde{F}({\\bf x}) = G({\\bf x})$.\n\\item For ${\\bf x}$ ``sufficiently far away'' from $\\bar{{\\bf x}}$, $\\tilde{F}({\\bf x}) \\simeq F({\\bf x})$.\n\\item The function $\\tilde{F}({\\bf x})$ is ``approximately\" smooth submodular.\n\\end{itemize}\nOnce we have $\\tilde{F}({\\bf x})$, we can fix it to obtain a smooth submodular\nfunction $\\hat{F}({\\bf x})$, which is still close to the original function $F({\\bf x})$.\nWe also fix $G({\\bf x})$ in the same way, to obtain a function $\\hat{G}({\\bf x})$\nwhich is equal to $\\hat{F}({\\bf x})$ whenever ${\\bf x}$ is close to $\\bar{{\\bf x}}$. We defer\nthis step until the end.\n\nWe define $\\tilde{F}({\\bf x})$ as a convex linear combination\nof $F({\\bf x})$ and $G({\\bf x})$, guided by a ``smooth transition'' function, depending\non the distance of ${\\bf x}$ from $\\bar{{\\bf x}}$. The form that we use is the following:\\footnote{We remark\nthat a construction analogous to \\cite{MSV08}\nwould be $\\tilde{F}({\\bf x}) = F({\\bf x}) - \\phi(H({\\bf x}))$ where $H({\\bf x}) = F({\\bf x}) - G({\\bf x})$.\nWhile this makes the analysis easier in \\cite{MSV08}, it cannot be used in general.\nRoughly speaking, the problem is that in general the partial derivatives of $H({\\bf x})$\nare not bounded in any way by the value of $H({\\bf x})$.}\n$$ \\tilde{F}({\\bf x}) = (1-\\phi(D({\\bf x}))) F({\\bf x}) + \\phi(D({\\bf x})) G({\\bf x}) $$\nwhere $\\phi:{\\boldmath R}_+ \\rightarrow [0,1]$ is a suitable smooth function,\nand\n$$ D({\\bf x}) = ||{\\bf x} - \\bar{{\\bf x}}||^2 = \\sum_i (x_i - \\bar{x}_i)^2.$$\nThe idea is that when ${\\bf x}$ is close to $\\bar{{\\bf x}}$, $\\phi(D({\\bf x}))$ should be close to $1$,\ni.e. the convex linear combination should give most of the weight to $G({\\bf x})$.\nThe weight should shift gradually to $F({\\bf x})$ as ${\\bf x}$ gets further away from $\\bar{{\\bf x}}$.\nTherefore, we define $\\phi(t) = 1$ in a small interval $t \\in [0,\\delta]$,\nand $\\phi(t)$ tends to $0$ as $t$ increases.\nThis guarantees that $\\tilde{F}({\\bf x}) = G({\\bf x})$ whenever\n$D({\\bf x}) = ||{\\bf x} - \\bar{{\\bf x}}||^2 \\leq \\delta$.\nWe defer the precise construction of $\\phi(t)$ to Lemma~\\ref{lemma:phi-construction},\nafter we determine what properties we need from $\\phi(t)$.\nNote that regardless of the definition of $\\phi(t)$, $\\tilde{F}({\\bf x})$\nis symmetric with respect to ${\\cal G}$, since $F({\\bf x}), G({\\bf x})$ and $D({\\bf x})$ are.\n\n\\\n\n\\paragraph{Analysis of the construction}\nDue to the construction of $\\tilde{F}({\\bf x})$, it is clear that when\n$D({\\bf x}) = ||{\\bf x}-\\bar{{\\bf x}}||^2$ is small, $\\tilde{F}({\\bf x}) = G({\\bf x})$.\nWhen $D({\\bf x})$ is large, $\\tilde{F}({\\bf x}) \\simeq F({\\bf x})$.\nThe main issue, however, is whether we can say something about the first\nand second partial derivatives of $\\tilde{F}$. This is crucial for\nthe properties of monotonicity and submodularity, which we would like to preserve.\nLet us write $\\tilde{F}({\\bf x})$ as\n$$ \\tilde{F}({\\bf x}) = F({\\bf x}) - \\phi(D({\\bf x})) H({\\bf x}) $$\nwhere $H({\\bf x}) = F({\\bf x}) - G({\\bf x})$. By differentiating once, we get\n\\begin{equation}\n\\label{eq:partdiff}\n\\partdiff{\\tilde{F}}{x_i} = \\partdiff{F}{x_i} - \\phi(D({\\bf x})) \\partdiff{H}{x_i}\n - \\phi'(D({\\bf x})) \\partdiff{D}{x_i} H({\\bf x})\n\\end{equation}\nand by differentiating twice,\n\\begin{eqnarray}\n\\label{eq:mixdiff}\n \\mixdiff{\\tilde{F}}{x_i}{x_j} & = & \\mixdiff{F}{x_i}{x_j}\n- \\phi({D}({\\bf x})) \\mixdiff{H}{x_i}{x_j}\n - \\phi''({D}({\\bf x})) \\partdiff{{D}}{x_i} \\partdiff{{D}}{x_j} H({\\bf x}) \\\\\n & & - \\phi'({D}({\\bf x})) \\left( \\partdiff{{D}}{x_j} \\partdiff{H}{x_i}\n + \\mixdiff{{D}}{x_i}{x_j} H({\\bf x}) + \\partdiff{{D}}{x_i} \\partdiff{H}{x_j} \\right) \\nonumber.\n\\end{eqnarray}\n\nThe first two terms on the right-hand sides of (\\ref{eq:partdiff}) and\n(\\ref{eq:mixdiff}) are not bothering us, because they form\nconvex linear combinations of the derivatives of $F({\\bf x})$ and $G({\\bf x})$,\nwhich have the properties that we need. The remaining terms might cause\nproblems, however, and we need to estimate them.\n\nOur strategy is to define $\\phi(t)$ in such a way that it eliminates\nthe influence of partial derivatives of $D$ and $H$ where they become too large.\nRoughly speaking, $D$ and $H$ have negligible partial derivatives when ${\\bf x}$\nis very close to $\\bar{{\\bf x}}$. As ${\\bf x}$ moves away from $\\bar{{\\bf x}}$, the partial\nderivatives grow but then the behavior of $\\phi(t)$ must be such that\ntheir influence is supressed.\n\nWe start with the following important claim.\\footnote{\nWe remind the reader that $\\nabla F$, the gradient of $F$, is a vector\nwhose coordinates are the first partial derivatives $\\partdiff{F}{x_i}$.\nWe denote by $\\nabla F |_{{\\bf x}}$ the gradient evaluated at ${\\bf x}$.}\n\n\\begin{lemma}\n\\label{lemma:grad-symmetry}\nAssume that $F:[0,1]^X \\rightarrow {\\boldmath R}$ is differentiable and invariant under a group of\npermutations of coordinates ${\\cal G}$. Let $G({\\bf x}) = F(\\bar{{\\bf x}})$,\nwhere $\\bar{{\\bf x}} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x})]$.\nThen for any ${\\bf x} \\in [0,1]^X$,\n $$ \\nabla{G}|_{\\bf x} = \\nabla{F}|_{\\bar{{\\bf x}}}.$$\n\\end{lemma}\n\n\\begin{proof}\nTo avoid confusion, we use ${\\bf x}$ for the arguments of the functions $F$ and $G$,\nand ${\\bf u}$, $\\bar{{\\bf u}}$, etc. for points where their partial derivatives are evaluated.\nTo rephrase, we want to prove that for any point ${\\bf u}$ and any coordinate $i$,\nthe partial derivatives of $F$ and $G$ evaluated at $\\bar{u}$ are equal:\n$\\partdiff{G}{x_i} \\Big|_{{\\bf x}={\\bf u}} = \\partdiff{F}{x_i} \\Big|_{{\\bf x}=\\bar{{\\bf u}}}$.\n\nFirst, consider $F({\\bf x})$. We assume that $F({\\bf x})$ is invariant under\na group of permutations of coordinates $\\cal G$,\ni.e. $F({\\bf x}) = F(\\sigma({\\bf x}))$ for any $\\sigma \\in {\\cal G}$.\nDifferentiating both sides at ${\\bf x}={\\bf u}$, we get by the chain rule:\n$$ \\partdiff{F}{x_i} \\Big|_{{\\bf x}={\\bf u}} =\n \\sum_j \\partdiff{F}{x_j} \\Big|_{{\\bf x}=\\sigma({\\bf u})} \\partdiff{}{x_i} (\\sigma({\\bf x}))_j\n = \\sum_j \\partdiff{F}{x_j} \\Big|_{{\\bf x}=\\sigma({\\bf u})} \\partdiff{x_{\\sigma(j)}}{x_i}. $$\nHere, $\\partdiff{x_{\\sigma(j)}}{x_i} = 1$ if $\\sigma(j) = i$, and $0$ otherwise.\nTherefore,\n$$ \\partdiff{F}{x_i} \\Big|_{{\\bf x}={\\bf u}} =\n \\partdiff{F}{x_{\\sigma^{-1}(i)}} \\Big|_{{\\bf x}=\\sigma({\\bf u})}. $$\nNow, if we evaluate the left-hand side at $\\bar{{\\bf u}}$, the right-hand side is evaluated\nat $\\sigma(\\bar{{\\bf u}}) = \\bar{{\\bf u}}$, and hence for any $i$ and any $\\sigma \\in {\\cal G}$,\n\\begin{equation}\n\\label{eq:F-inv}\n\\partdiff{F}{x_i} \\Big|_{{\\bf x}=\\bar{{\\bf u}}} = \\partdiff{F}{x_{\\sigma^{-1}(i)}} \\Big|_{{\\bf x}=\\bar{{\\bf u}}}.\n\\end{equation}\nTurning to $G({\\bf x}) = F(\\bar{{\\bf x}})$, let us write $\\partdiff{G}{x_i}$\nusing the chain rule:\n$$ \\partdiff{G}{x_i} \\Big|_{{\\bf x}={\\bf u}} = \\partdiff{}{x_i} F(\\bar{{\\bf x}}) \\Big|_{{\\bf x}={\\bf u}}\n = \\sum_j \\partdiff{F}{x_j} \\Big|_{{\\bf x}=\\bar{{\\bf u}}} \\cdot \\partdiff{\\bar{x}_j}{x_i}. $$\nWe have $\\bar{x}_j = {\\bf E}_{\\sigma \\in {\\cal G}}[{\\bf x}_{\\sigma(j)}]$, and\nso\n$$ \\partdiff{G}{x_i} \\Big|_{{\\bf x}={\\bf u}} = \\sum_j \\partdiff{F}{x_j} \\Big|_{{\\bf x}=\\bar{{\\bf u}}}\n \\cdot \\partdiff{}{x_i} {\\bf E}_{\\sigma \\in {\\cal G}}[x_{\\sigma(j)}] \n = {\\bf E}_{\\sigma \\in {\\cal G}} \\left[\\sum_j \\partdiff{F}{x_j} \\Big|_{x=\\bar{u}}\n \\cdot \\partdiff{x_{\\sigma(j)}}{x_i} \\right]. $$\nAgain, $\\partdiff{x_{\\sigma(j)}}{x_i} = 1$ if $\\sigma(j) = i$ and $0$ otherwise.\nConsequently, we obtain\n$$ \\partdiff{G}{x_i} \\Big|_{{\\bf x}={\\bf u}} = {\\bf E}_{\\sigma \\in {\\cal G}}\n \\left[ \\partdiff{F}{x_{\\sigma^{-1}(i)}} \\Big|_{{\\bf x}=\\bar{{\\bf u}}} \\right]\n = \\partdiff{F}{x_i} \\Big|_{{\\bf x}=\\bar{{\\bf u}}} $$\nwhere we used Eq. (\\ref{eq:F-inv}) to remove the dependence on $\\sigma \\in {\\cal G}$.\n\\end{proof}\n\n\nObserve that the symmetrization operation $\\bar{{\\bf x}}$ is idempotent, i.e.\n$\\bar{\\bar{{\\bf x}}} = \\bar{{\\bf x}}$.\nBecause of this, we also get $\\nabla{G}|_{\\bar{{\\bf x}}} = \\nabla{F}|_{\\bar{{\\bf x}}}$.\nNote that $G(\\bar{{\\bf x}}) = F(\\bar{{\\bf x}})$ follows from the definition,\nbut it is not obvious that the same holds for gradients, since their\ndefinition involves points where $G({\\bf x}) \\neq F({\\bf x})$. For second partial\nderivatives, the equality no longer holds, as can be seen from\na simple example such as $F(x_1,x_2) = 1 - (1-x_1)(1-x_2)$,\n$G(x_1,x_2) = 1 - (1-\\frac{x_1+x_2}{2})^2$.\n\nNext, we show that the functions $F({\\bf x})$ and $G({\\bf x})$ are very similar\nin the close vicinity of the region where $\\bar{{\\bf x}} = {\\bf x}$. Recall our definitions:\n$H({\\bf x}) = F({\\bf x}) - G({\\bf x})$, $D({\\bf x}) = ||{\\bf x} - \\bar{{\\bf x}}||^2$.\nBased on Lemma~\\ref{lemma:grad-symmetry}, we know that $H(\\bar{{\\bf x}}) = 0$\nand $\\nabla H|_{\\bar{{\\bf x}}} = 0$. In the following lemmas,\nwe present bounds on $H({\\bf x})$, $D({\\bf x})$ and their partial derivatives.\n\n\n\\begin{lemma}\n\\label{lemma:H-bounds}\nLet $f:2^X \\rightarrow [0,M]$ be invariant under a permutation group $\\cal G$.\nLet $\\bar{{\\bf x}} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x})]$,\n$D({\\bf x}) = ||{\\bf x} - \\bar{{\\bf x}}||^2$ and $H({\\bf x}) = F({\\bf x}) - G({\\bf x})$ where $F({\\bf x}) = {\\bf E}[f(\\hat{{\\bf x}})]$\nand $G({\\bf x}) = F(\\bar{{\\bf x}})$. Then\n\\begin{enumerate}\n\\item $ |\\mixdiff{H}{x_i}{x_j}| \\leq 8M $ everywhere, for all $i,j$;\n\\item $ ||\\nabla H({\\bf x})|| \\leq 8M|X| \\sqrt{D({\\bf x})}$;\n\\item $ |H({\\bf x})| \\leq 8M|X| \\cdot D({\\bf x}). $\n\\end{enumerate}\n\\end{lemma}\n\n\\begin{proof}\nFirst, let us get a bound on the second partial derivatives.\nAssuming without loss of generality $x_i=x_j=0$, we have\\footnote{${\\bf x} \\vee {\\bf y}$\ndenotes the coordinate-wise maximum, $({\\bf x} \\vee {\\bf y})_i = \\max \\{x_i,y_i\\}$ and\n${\\bf x} \\wedge {\\bf y}$ denotes the coordinate-wise minimum, $({\\bf x} \\wedge {\\bf y})_i = \\min \\{x_i,y_i\\}$.}\n$$ \\mixdiff{F}{x_i}{x_j} = {\\bf E}[f(\\hat{{\\bf x}} \\vee ({\\bf e}_i+{\\bf e}_j)) -\n f(\\hat{{\\bf x}} \\vee {\\bf e}_i) - f(\\hat{{\\bf x}} \\vee {\\bf e}_j) + f(\\hat{{\\bf x}})] $$\n(see \\cite{Vondrak08}). Consequently,\n$$ \\Big| \\mixdiff{F}{x_i}{x_j} \\Big| \\leq 4 \\max |f(S)| = 4 M.$$\nIt is a little bit more involved to analyze $\\mixdiff{G}{x_i}{x_j}$.\nSince $G({\\bf x}) = F(\\bar{{\\bf x}})$ and $\\bar{{\\bf x}} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x})]$,\nwe get by the chain rule:\n$$ \\mixdiff{G}{x_i}{x_j} = \\sum_{k,\\ell} \\mixdiff{F}{x_k}{x_\\ell}\n \\partdiff{\\bar{x}_k}{x_i} \\partdiff{\\bar{x}_\\ell}{x_j} \n = {\\bf E}_{\\sigma,\\tau \\in {\\cal G}} \\left[ \\sum_{k,\\ell} \\mixdiff{F}{x_k}{x_\\ell}\n \\partdiff{x_{\\sigma(k)}}{x_i} \\partdiff{x_{\\tau(\\ell)}}{x_j} \\right].$$\nIt is useful here to use the Kronecker symbol, $\\delta_{i,j}$,\nwhich is $1$ if $i=j$ and $0$ otherwise. Note that\n$\\partdiff{x_{\\sigma(k)}}{x_i} = \\delta_{i,\\sigma(k)} = \\delta_{\\sigma^{-1}(i),k}$,\netc. Using this notation, we get\n$$ \\mixdiff{G}{x_i}{x_j} = {\\bf E}_{\\sigma,\\tau \\in {\\cal G}} \\left[ \\sum_{k,\\ell}\n\\mixdiff{F}{x_k}{x_\\ell} \\delta_{\\sigma^{-1}(i),k} \\delta_{\\sigma^{-1}(j),\\ell} \\right] \n = {\\bf E}_{\\sigma,\\tau \\in {\\cal G}}\n \\left[ \\mixdiff{F}{x_{\\sigma^{-1}(i)}}{x_{\\tau^{-1}(j)}} \\right], $$\n$$ \\Big| \\mixdiff{G}{x_i}{x_j} \\Big| \\leq {\\bf E}_{\\sigma,\\tau \\in {\\cal G}}\n\\left[ \\Big| \\mixdiff{F}{x_{\\sigma^{-1}(i)}}{x_{\\tau^{-1}(j)}} \\Big| \\right]\n \\leq 4 M $$\nand therefore\n$$ \\Big| \\mixdiff{H}{x_i}{x_j} \\Big|\n = \\Big| \\mixdiff{F}{x_i}{x_j} - \\mixdiff{G}{x_i}{x_j} \\Big|\n \\leq 8 M.$$\n\nNext, we estimate $\\partdiff{H}{x_i}$ at a given point ${\\bf u}$, depending\non its distance from $\\bar{{\\bf u}}$. Consider\nthe line segment between $\\bar{{\\bf u}}$ and ${\\bf u}$. The function $H({\\bf x}) = F({\\bf x}) - G({\\bf x})$\nis $C_\\infty$-differentiable, and hence we can apply the mean value theorem\nto $\\partdiff{H}{x_i}$: There exists a point $\\tilde{{\\bf u}}$ on the line segment\n$[\\bar{{\\bf u}}, {\\bf u}]$ such that\n$$ \\partdiff{H}{x_i} \\Big|_{{\\bf x}={\\bf u}} - \\partdiff{H}{x_i} \\Big|_{{\\bf x}=\\bar{{\\bf u}}}\n = \\sum_j \\mixdiff{H}{x_j}{x_i} \\Big|_{{\\bf x}=\\tilde{{\\bf u}}} (u_j-\\bar{u}_j).$$\nRecall that $\\partdiff{H}{x_i} \\Big|_{{\\bf x}=\\bar{{\\bf u}}} = 0$.\nApplying the Cauchy-Schwartz inequality to the right-hand side, we get\n$$ \\left( \\partdiff{H}{x_i} \\Big|_{{\\bf x}={\\bf u}} \\right)^2\n \\leq \\sum_j \\left( \\mixdiff{H}{x_j}{x_i} \\Big|_{{\\bf x}=\\tilde{{\\bf u}}} \\right)^2\n || {\\bf u} - \\bar{{\\bf u}} ||^2 \\leq (8M)^2 |X| ||{\\bf u} - \\bar{{\\bf u}}||^2.$$\nAdding up over all $i \\in X$, we obtain\n$$ || \\nabla H({\\bf u}) ||^2 = \\sum_{i}\\left( \\partdiff{H}{x_i} \\Big|_{{\\bf x}={\\bf u}} \\right)^2\n \\leq (8M|X|)^2 ||{\\bf u} - \\bar{{\\bf u}}||^2.$$\nFinally, we estimate the growth of $H({\\bf u})$. Again, by the mean value theorem,\nthere is a point $\\tilde{{\\bf u}}$ on the line segment $[\\bar{{\\bf u}},{\\bf u}]$, such that\n$$ H({\\bf u}) - H(\\bar{{\\bf u}}) = ({\\bf u} - \\bar{{\\bf u}}) \\cdot \\nabla H(\\tilde{{\\bf u}}).$$\nUsing $H(\\bar{{\\bf u}}) = 0$, the Cauchy-Schwartz inequality and the above bound\non $\\nabla H$,\n$$ (H({\\bf u}))^2 \\leq ||\\nabla H({\\tilde{{\\bf u}}})||^2 ||{\\bf u} - \\bar{{\\bf u}}||^2\n \\leq (8M|X|)^2 ||\\tilde{{\\bf u}}-\\bar{{\\bf u}}||^2 ||{\\bf u}-\\bar{{\\bf u}}||^2. $$\nClearly, $||\\tilde{{\\bf u}} - \\bar{{\\bf u}}|| \\leq ||{\\bf u} - \\bar{{\\bf u}}||$, and therefore\n$$ |H({\\bf u})| \\leq 8M|X| \\cdot ||{\\bf u}-\\bar{{\\bf u}}||^2.$$\n\\end{proof}\n\n\\begin{lemma}\n\\label{lemma:D-bounds}\nFor the function $D({\\bf x}) = ||{\\bf x} - \\bar{{\\bf x}}||^2$, we have\n\\begin{enumerate}\n\\item $\\nabla D = 2({\\bf x} - \\bar{{\\bf x}})$, and therefore\n $||\\nabla D|| = 2 \\sqrt{D({\\bf x})}$.\n\\item For all $i,j$, $|\\mixdiff{D}{x_i}{x_j}| \\leq 2$.\n\\end{enumerate}\n\\end{lemma}\n\n\\begin{proof}\nLet us write $D({\\bf x})$ as\n$$ D({\\bf x}) = \\sum_i (x_i - \\bar{x}_i)^2 = \\sum_i {\\bf E}_{\\sigma \\in {\\cal G}}[x_i - x_{\\sigma(i)}]\n {\\bf E}_{\\tau \\in {\\cal G}}[x_i - x_{\\tau(i)}]. $$\nTaking the first partial derivative,\n$$ \\partdiff{D}{x_j} = 2 \\sum_i {\\bf E}_{\\sigma \\in {\\cal G}}[x_i - x_{\\sigma(i)}]\n \\partdiff{}{x_j} {\\bf E}_{\\tau \\in {\\cal G}}[x_i - x_{\\tau(i)}].$$\nAs before, we have $\\partdiff{x_i}{x_j} = \\delta_{ij}$.\nUsing this notation, we get\n\\begin{eqnarray*}\n\\partdiff{D}{x_j} & = & 2 \\sum_i {\\bf E}_{\\sigma \\in {\\cal G}}[x_i - x_{\\sigma(i)}]\n {\\bf E}_{\\tau \\in {\\cal G}}[\\delta_{ij} - \\delta_{\\tau(i),j}] \\\\\n & = & 2 \\sum_i {\\bf E}_{\\sigma,\\tau \\in {\\cal G}}[(x_i - x_{\\sigma(i)})\n (\\delta_{ij} - \\delta_{i,\\tau^{-1}(j)})] \\\\\n & = & 2 \\ {\\bf E}_{\\sigma,\\tau \\in {\\cal G}}[x_j - x_{\\sigma(j)} - x_{\\tau^{-1}(j)}\n + x_{\\sigma(\\tau^{-1}(j))}].\n\\end{eqnarray*}\nSince the distributions of $\\sigma(j)$, $\\tau^{-1}(j)$ and $\\sigma(\\tau^{-1}(j))$\nare the same, we obtain\n$$ \\partdiff{D}{x_j} = 2 \\ {\\bf E}_{\\sigma \\in {\\cal G}}[x_j - x_{\\sigma(j)}]\n = 2(x_j - \\bar{x}_j) $$\nand\n$$ ||\\nabla D||^2 = \\sum_j \\Big| \\partdiff{D}{x_j} \\Big|^2\n = 4 \\sum_j (x_j - \\bar{x}_j)^2 = 4 D({\\bf x}).$$\n\nFinally, the second partial derivatives are\n$$ \\mixdiff{D}{x_i}{x_j} = 2 \\partdiff{}{x_i} (x_j - \\bar{x}_j)\n = 2 \\partdiff{}{x_i} {\\bf E}_{\\sigma \\in {\\cal G}}[x_j - x_{\\sigma(j)}]\n = 2 {\\bf E}_{\\sigma \\in {\\cal G}}[\\delta_{ij} - \\delta_{i,\\sigma(j)}] $$\nwhich is clearly bounded by $2$ in the absolute value.\n\\end{proof}\n\nNow we come back to $\\tilde{F}({\\bf x})$ and its partial derivatives.\nRecall equations (\\ref{eq:partdiff}) and (\\ref{eq:mixdiff}).\nThe problematic terms are those involving $\\phi'(D({\\bf x}))$ and $\\phi''(D({\\bf x}))$.\nUsing our bounds on $H({\\bf x})$, $D({\\bf x})$ and their derivatives, however,\nwe notice that $\\phi'(D({\\bf x}))$ always appears with factors on the order\nof $D({\\bf x})$ and $\\phi''(D({\\bf x}))$ appears with factors on the order of $(D({\\bf x}))^2$.\nThus, it is sufficient if $\\phi(t)$ is defined so that we have control\nover $t \\phi'(t)$ and $t^2 \\phi''(t)$. The following lemma describes\nthe function that we need.\n\n\\begin{lemma}\n\\label{lemma:phi-construction}\nFor any $\\alpha, \\beta > 0$, there is $\\delta \\in (0, \\beta)$\nand a function $\\phi:{\\boldmath R}_+ \\rightarrow [0,1]$ with an absolutely continuous first derivative\nsuch that\n\\begin{enumerate}\n\\item For $t \\leq \\delta$, $\\phi(t) = 1$.\n\\item For $t \\geq \\beta$, $\\phi(t) < e^{-1\/\\alpha}$.\n\\item For all $t \\geq 0$, $|t \\phi'(t)| \\leq 4 \\alpha$.\n\\item For almost all $t \\geq 0$, $|t^2 \\phi''(t)| \\leq 10 \\alpha$.\n\\end{enumerate}\n\\end{lemma}\n\n\\begin{proof}\nFirst, observe that if we prove the lemma for some particular value $\\beta>0$, we can also prove prove it\nfor any other value $\\tilde{\\beta}>0$, by modifying the function as follows: $\\tilde{\\phi}(t) = \\phi(\\beta t \/ \\beta')$.\nThis corresponds to a scaling of the parameter $t$ by $\\beta' \/ \\beta$. Observe that then $|t \\tilde{\\phi}'(t)| = |\\frac{\\beta}{\\beta'} t \\phi'(\\beta t \/ \\beta')| \\leq 4 \\alpha$ and $|t^2 \\tilde{\\phi}''(t)| = |(\\frac{\\beta}{\\beta'})^2 t^2 \\phi''(\\beta t \/ \\beta')| \\leq 10 \\alpha$, so the conditions are still satisfied.\n\nTherefore, we can assume without\nloss of generality that $\\beta>0$ is a value of our choice,\nfor example $\\beta = e^{1\/(2\\alpha^2)}+1$.\nIf we want to prove the result for a different value of $\\beta$,\nwe can just scale the argument $t$ and the constant $\\delta$\nby $\\beta \/ (e^{1\/(2\\alpha^2)}+1)$; the bounds on\n$t \\phi'(t)$ and $t^2 \\phi''(t)$ still hold.\n\nWe can assume that $\\alpha \\in (0,\\frac18)$ because for larger $\\alpha$,\nthe statement only gets weaker. As we argued, we can assume WLOG that\n$\\beta = e^{1\/(2\\alpha^2)}+1$. We set $\\delta = 1$ and\n$\\delta_2 = 1 + (1+\\alpha)^{-1\/2} \\leq 2$.\n(We remind the reader that in general these values will be scaled depending\non the actual value of $\\beta$.) We define the function as follows:\n\\begin{enumerate}\n\\item $\\phi(t) = 1$ for $t \\in [0,\\delta]$.\n\\item $\\phi(t) = 1 - \\alpha (t-1)^2$ for $t \\in [\\delta, \\delta_2]$.\n\\item $\\phi(t) = (1+\\alpha)^{-1-\\alpha} (t - 1)^{-2 \\alpha}$\n for $t \\in [\\delta_2, \\infty)$.\n\\end{enumerate}\n\n\\noindent\nLet's verify the properties of $\\phi(t)$.\nFor $t \\in [0,\\delta]$, we have\n$\\phi'(t) = \\phi''(t) = 0$. For $t \\in [\\delta,\\delta_2]$, we have\n$$ \\phi'(t) = -2\\alpha \\left(t-1 \\right), \\ \\ \\ \\ \\ \\ \n \\phi''(t) = -2\\alpha, $$\nand for $t \\in [\\delta_2, \\infty)$,\n$$ \\phi'(t) = -{2\\alpha}{(1+\\alpha)^{-1-\\alpha}} \\left( t-1 \\right)^{-2\\alpha-1}, $$\n$$ \\phi''(t) = {2\\alpha(1+2\\alpha)}{(1+\\alpha)^{-1-\\alpha}} \\left( t-1 \\right)^{-2\\alpha-2}. $$\nFirst, we check that the values and first derivatives agree at the breakpoints.\nFor $t = \\delta = 1$, we get $\\phi(1) = 1$ and $\\phi'(1) = 0$.\nFor $t = \\delta_2 = 1 + (1+\\alpha)^{-1\/2}$, we get\n$\\phi(\\delta_2) = (1+\\alpha)^{-1}$ and\n$\\phi'(\\delta_2) = -2 \\alpha (1+\\alpha)^{-1\/2}$.\nNext, we need to check is that $\\phi(t)$ is very small\nfor $t \\geq \\beta$. The function is decreasing for $t > \\beta$,\ntherefore it is enough to check $t = \\beta = e^{1\/(2\\alpha^2)}+1$:\n$$ \\phi\\left(\\beta \\right) = (1+\\alpha)^{-1-\\alpha} (\\beta - 1)^{-2\\alpha}\n \\leq (\\beta-1)^{-2\\alpha} = e^{-1\/\\alpha}. $$\nThe derivative bounds are satisfied trivially for $t \\in [0,\\delta]$.\nFor $t \\in [\\delta, \\delta_2]$, using $t \\leq \\delta_2 = 1 + (1+\\alpha)^{-1\/2}$,\n$$ |t \\phi'(t)| = t \\cdot 2 \\alpha (t-1)\n \\leq 2 \\alpha (1 + (1+\\alpha)^{-1\/2}) (1+\\alpha)^{-1\/2} \\leq 4 \\alpha $$\nand using $\\alpha \\in (0,\\frac18)$,\n$$ |t^2 \\phi''(t)| = t^2 \\cdot 2 \\alpha \\leq 2 \\alpha (1 + (1+\\alpha)^{-1\/2})^2\n \\leq 8 \\alpha.$$\nFor $t \\in [\\delta_2,\\infty)$, using $t-1 \\geq (1+\\alpha)^{-1\/2}$,\n\\begin{eqnarray*}\n|t \\phi'(t)| & = & t \\cdot \\frac{2 \\alpha}{(1+\\alpha)^{1+\\alpha}} \\left(t-1 \\right)^{-2\\alpha-1} \n = \\frac{2 \\alpha}{1+\\alpha} \\left( (1 + \\alpha)\n\\left( t-1 \\right)^2 \\right)^{-\\alpha} \\frac{t}{t-1} \\\\\n& \\leq & \\frac{2 \\alpha}{1 + \\alpha} \\cdot \\frac{t}{t-1}\n \\leq \\frac{2 \\alpha}{1+\\alpha} \\cdot \\frac{1 + (1+\\alpha)^{-1\/2}}{(1+\\alpha)^{-1\/2}}\n = 2 \\alpha \\cdot \\frac{1 + (1+\\alpha)^{-1\/2}}{(1+\\alpha)^{1\/2}} \\leq 4 \\alpha\n\\end{eqnarray*}\nand finally, using $\\alpha \\in (0,\\frac18)$,\n\\begin{eqnarray*}\n |t^2 \\phi''(t)| & = & t^2 \\cdot \\frac{2\\alpha(1+2\\alpha)}{(1+\\alpha)^{1+\\alpha}}\n \\left( t-1 \\right)^{-2\\alpha-2} \n = \\frac{2\\alpha(1+2\\alpha)}{1+\\alpha}\n \\left( (1 + \\alpha) \\left( t-1 \\right)^2 \\right)^{-\\alpha}\n \\left( \\frac{t}{t-1} \\right)^2 \\\\\n& \\leq & \\frac{2\\alpha(1+2\\alpha)}{1+\\alpha} \\cdot\n \\left( \\frac{1 + (1+\\alpha)^{-1\/2}}{(1+\\alpha)^{-1\/2}} \\right)^2\n \\leq 8 \\alpha (1 + 2 \\alpha) \\leq 10 \\alpha.\n\\end{eqnarray*}\n\\end{proof}\n\nUsing Lemmas~\\ref{lemma:H-bounds}, \\ref{lemma:D-bounds} and \\ref{lemma:phi-construction},\nwe now prove bounds on the derivatives of $\\tilde{F}({\\bf x})$.\n\n\\begin{lemma}\n\\label{lemma:F-bounds}\nLet $ \\tilde{F}({\\bf x}) = (1-\\phi(D({\\bf x}))) F({\\bf x}) + \\phi(D({\\bf x})) G({\\bf x}) $\nwhere $F({\\bf x}) = {\\bf E}[f(\\hat{{\\bf x}})]$, $f:2^X \\rightarrow [0,M]$,\n$G({\\bf x}) = F(\\bar{{\\bf x}})$, $D({\\bf x})=||{\\bf x}-\\bar{{\\bf x}}||^2$\nare as above and $\\phi(t)$ is provided by Lemma~\\ref{lemma:phi-construction}\nfor a given $\\alpha>0$.\nThen, whenever $\\mixdiff{F}{x_i}{x_j} \\leq 0$,\n$$\n \\mixdiff{\\tilde{F}}{x_i}{x_j}\n \\leq 512 M |X| \\alpha. $$\nIf, in addition, $\\partdiff{F}{x_i} \\geq 0$, then\n$$\n\\partdiff{\\tilde{F}}{x_i}\n \\geq -64 M |X| \\alpha. $$\n\\end{lemma}\n\n\\begin{proof}\nWe have\n$ \\tilde{F}({\\bf x}) = F({\\bf x}) - \\phi(D({\\bf x})) H({\\bf x}) $\nwhere $H({\\bf x}) = F({\\bf x}) - G({\\bf x})$. By differentiating once, we get\n$$ \\partdiff{\\tilde{F}}{x_i} = \\partdiff{F}{x_i}\n - \\phi(D({\\bf x})) \\partdiff{H}{x_i} - \\phi'(D({\\bf x})) \\partdiff{D}{x_i} H({\\bf x}), $$\ni.e.\n$$ \\Big| \\partdiff{\\tilde{F}}{x_i} -\n \\left(\\partdiff{F}{x_i} - \\phi(D({\\bf x})) \\partdiff{H}{x_i} \\right) \\Big|\n = \\Big| \\phi'(D({\\bf x})) \\partdiff{D}{x_i} H({\\bf x}) \\Big|.$$\nBy Lemma~\\ref{lemma:H-bounds} and \\ref{lemma:D-bounds},\nwe have $|\\partdiff{D}{x_i}| = 2 |x_i - \\bar{x}_i| \\leq 2$\nand $|H({\\bf x})| \\leq 8M|X| \\cdot D({\\bf x})$. Therefore,\n$$ \\Big| \\partdiff{\\tilde{F}}{x_i} -\n \\left(\\partdiff{F}{x_i} - \\phi(D({\\bf x})) \\partdiff{H}{x_i} \\right) \\Big|\n \\leq 16 M |X| D({\\bf x}) \\cdot \\Big| \\phi'(D({\\bf x})) \\Big|.$$\nBy Lemma~\\ref{lemma:phi-construction}, $|D({\\bf x}) \\phi'(D({\\bf x}))| \\leq 4 \\alpha$, and hence\n$$ \\Big| \\partdiff{\\tilde{F}}{x_i} -\n \\left(\\partdiff{F}{x_i} - \\phi(D({\\bf x})) \\partdiff{H}{x_i} \\right) \\Big|\n \\leq 64 M |X| \\alpha.$$\nAssuming that $\\partdiff{F}{x_i} \\geq 0$, we also have $\\partdiff{G}{x_i} \\geq 0$\n(see Lemma~\\ref{lemma:grad-symmetry}) and therefore,\n$ \\partdiff{F}{x_i} - \\phi(D({\\bf x})) \\partdiff{H}{x_i}\n = (1-\\phi(D({\\bf x}))) \\partdiff{F}{x_i} + \\phi(D({\\bf x})) \\partdiff{G}{x_i} \\geq 0$.\nConsequently,\n$$ \\partdiff{\\tilde{F}}{x_i} \\geq - 64 M |X| \\alpha.$$\nBy differentiating $\\tilde{F}$ twice, we obtain\n\\begin{eqnarray*}\n\\mixdiff{\\tilde{F}}{x_i}{x_j} & = &\n \\mixdiff{F}{x_i}{x_j} - \\phi({D}({\\bf x})) \\mixdiff{H}{x_i}{x_j}\n - \\phi''({D}({\\bf x})) \\partdiff{{D}}{x_i} \\partdiff{{D}}{x_j} H({\\bf x}) \\\\\n& & - \\phi'({D}({\\bf x})) \\left( \\partdiff{{D}}{x_j} \\partdiff{H}{x_i}\n + \\mixdiff{{D}}{x_i}{x_j} H({\\bf x})\n + \\partdiff{{D}}{x_i} \\partdiff{H}{x_j} \\right).\n\\end{eqnarray*}\nAgain, we use Lemma~\\ref{lemma:H-bounds} and \\ref{lemma:D-bounds}\nto bound $|H({\\bf x})| \\leq 8M|X| D({\\bf x})$, $|\\partdiff{H}{x_i}| \\leq 8M|X| \\sqrt{D({\\bf x})}$,\n $|\\mixdiff{H}{x_i}{x_j}| \\leq 8M$, $|\\partdiff{D}{x_i}| \\leq 2 \\sqrt{D({\\bf x})}$\nand $|\\mixdiff{D}{x_i}{x_j}| \\leq 2$. We get\n\\begin{eqnarray*}\n\\Big| \\mixdiff{\\tilde{F}}{x_i}{x_j} - \\mixdiff{F}{x_i}{x_j} +\n \\phi({D}({\\bf x})) \\mixdiff{H}{x_i}{x_j} \\Big|\n \\leq 32 M|X| \\Big| D^2({\\bf x}) \\phi''({D}({\\bf x})) \\Big|\n + 48 M|X| \\Big| D({\\bf x}) \\phi'({D}({\\bf x})) \\Big|\n\\end{eqnarray*}\nObserve that $\\phi'(D({\\bf x}))$ appears with $D({\\bf x})$ and\n$\\phi''(D({\\bf x}))$ appears with $(D({\\bf x}))^2$.\nBy Lemma~\\ref{lemma:phi-construction}, $|D({\\bf x}) \\phi'(D({\\bf x}))| \\leq 4 \\alpha$\nand $|D^2({\\bf x}) \\phi''(D({\\bf x}))| \\leq 10 \\alpha$. Therefore,\n\\begin{eqnarray*}\n\\Big| \\mixdiff{\\tilde{F}}{x_i}{x_j} - \\left( \\mixdiff{F}{x_i}{x_j}\n - \\phi({D}({\\bf x})) \\mixdiff{H}{x_i}{x_j} \\right) \\Big|\n \\leq 320 M|X|\\alpha + 192 M|X| \\alpha = 512 M |X| \\alpha.\n\\end{eqnarray*}\nIf $\\mixdiff{F}{x_i}{x_j} \\leq 0$ for all $i,j$, then also $\\mixdiff{G}{x_i}{x_j} =\n {\\bf E}_{\\sigma,\\tau \\in {\\cal G}}[\\mixdiff{F}{x_{\\sigma^{-1}(i)}}{x_{\\tau^{-1}(j)}}] \\leq 0$\n (see the proof of Lemma~\\ref{lemma:H-bounds}). Also, \n$\\mixdiff{F}{x_i}{x_j} - \\phi({D}({\\bf x})) \\mixdiff{H}{x_i}{x_j}\n = \\phi(D({\\bf x})) \\mixdiff{F}{x_i}{x_j} + (1-\\phi(D({\\bf x}))) \\mixdiff{G}{x_i}{x_j} \\leq 0$.\nWe obtain\n\\begin{eqnarray*}\n\\mixdiff{\\tilde{F}}{x_i}{x_j} \\leq 512 M |X| \\alpha.\n\\end{eqnarray*}\n\\end{proof}\n\nFinally, we can finish the proof of Lemma~\\ref{lemma:final-fix}.\n\n\\begin{proof}[Proof of Lemma~\\ref{lemma:final-fix}]\nLet $\\epsilon>0$ and $f:2^X \\rightarrow [0,M]$.\nWe choose $\\beta = \\frac{\\epsilon}{16 M|X|}$, so that $|H({\\bf x})| = |F({\\bf x}) - G({\\bf x})| \\leq \\epsilon\/2$\nwhenever $D({\\bf x}) = ||{\\bf x} - \\bar{{\\bf x}}||^2 \\leq \\beta$\n(due to by Lemma~\\ref{lemma:H-bounds}, which states that $|H({\\bf x})| \\leq 8 M |X| D({\\bf x})$).\nAlso, let $\\alpha = \\frac{\\epsilon}{2000 M |X|^3}$.\nFor these values of $\\alpha,\\beta>0$,\nlet $\\delta > 0$ and $\\phi:{\\boldmath R}_+ \\rightarrow [0,1]$\nbe provided by Lemma~\\ref{lemma:phi-construction}.\nWe define\n$$ \\tilde{F}({\\bf x}) = (1-\\phi(D({\\bf x}))) F({\\bf x}) + \\phi(D({\\bf x})) G({\\bf x}). $$\nLemma~\\ref{lemma:F-bounds} provides bounds on the first and second\npartial derivatives of $\\tilde{F}({\\bf x})$.\nFinally, we modify $\\tilde{F}({\\bf x})$ so that it satisfies the required\nconditions (submodularity and optionally monotonicity).\nFor that purpose, we add a suitable multiple of the following function:\n$$ J({\\bf x}) = |X|^2 + 3|X| \\sum_{i \\in X} x_i - \\left(\\sum_{i \\in X} x_i \\right)^2.$$\nWe have $0 \\leq J({\\bf x}) \\leq 3|X|^2$,\n$\\partdiff{J}{x_i} = 3|X| - 2 \\sum_{i \\in X} x_i \\geq |X|$.\nFurther, $\\mixdiff{J}{x_i}{x_j} = -2$. Note also that $J(\\bar{{\\bf x}}) = J({\\bf x})$,\nsince $J({\\bf x})$ depends only on the sum of all coordinates $\\sum_{i \\in X} x_i$.\nTo make $\\tilde{F}({\\bf x})$ submodular and optionally monotone, we define:\n$$ \\hat{F}({\\bf x}) = \\tilde{F}({\\bf x}) + 256 M |X| \\alpha J({\\bf x}),$$\n$$ \\hat{G}({\\bf x}) = G({\\bf x}) + 256 M |X| \\alpha J({\\bf x}). $$\nWe verify the properties of $\\hat{F}({\\bf x})$ and $\\hat{G}({\\bf x})$:\n\\begin{enumerate}\n\\item\nFor any ${\\bf x} \\in P({\\cal F})$, we have\n\\begin{eqnarray*}\n\\hat{G}({\\bf x}) & = & G({\\bf x}) + 256 M |X| \\alpha J({\\bf x}) \\\\\n & = & F(\\bar{{\\bf x}}) + 256 M |X| \\alpha J(\\bar{{\\bf x}}) \\\\\n & = & \\hat{F}(\\bar{{\\bf x}}).\n\\end{eqnarray*}\n\n\\item When $D({\\bf x}) = ||{\\bf x} - \\bar{{\\bf x}}||^2 \\geq \\beta$,\nLemma~\\ref{lemma:phi-construction} guarantees\nthat $0 \\leq \\phi(D({\\bf x})) < e^{-1\/\\alpha} \\leq \\alpha$ and\n\\begin{eqnarray*}\n|\\hat{F}({\\bf x}) - F({\\bf x})| & \\leq & \\phi(D({\\bf x})) |G({\\bf x}) - F({\\bf x})| + 256 M |X| \\alpha J({\\bf x}) \\\\\n & \\leq & \\alpha M + 768 M |X|^3 \\alpha \\\\\n & \\leq & \\epsilon\n\\end{eqnarray*}\nusing $0 \\leq F({\\bf x}), G({\\bf x}) \\leq M$, $|X| \\leq J({\\bf x}) \\leq 3|X|^2$ and $\\alpha = \\frac{\\epsilon}{2000 M |X|^3}$.\n\nWhen $D({\\bf x}) = ||{\\bf x} - \\bar{{\\bf x}}||^2 < \\beta$, we chose the value of $\\beta$ so that\n$|G({\\bf x}) - F({\\bf x})| < \\epsilon\/2$ and so by the above,\n\\begin{eqnarray*}\n|\\hat{F}({\\bf x}) - F({\\bf x})| & \\leq & \\phi(D({\\bf x})) |G({\\bf x}) - F({\\bf x})| + 256 M |X| \\alpha J({\\bf x}) \\\\\n & \\leq & \\epsilon\/2 + 768 M |X|^3 \\alpha \\\\\n & \\leq & \\epsilon.\n\\end{eqnarray*}\n\n\\item Due to Lemma~\\ref{lemma:phi-construction}, $\\phi(t) = 1$ for $t \\in [0,\\delta]$.\nHence, whenever $D({\\bf x}) = ||{\\bf x} - \\bar{{\\bf x}}||^2 \\leq \\delta$, we have\n$ \\tilde{F}({\\bf x}) = G({\\bf x}) = F(\\bar{{\\bf x}})$, which depends only on $\\bar{{\\bf x}}$.\nAlso, we have $\\hat{F}({\\bf x}) = \\hat{G}({\\bf x}) = F(\\bar{{\\bf x}}) + 256 M |X| \\alpha J({\\bf x})$\nand again, $J({\\bf x})$ depends only on $\\bar{{\\bf x}}$ (in fact, only on the average\nof all coordinates of ${\\bf x}$). Therefore, $\\hat{F}({\\bf x})$ and $\\hat{G}({\\bf x})$\nin this case depend only on $\\bar{{\\bf x}}$.\n\n\\item The first partial derivatives of $\\hat{F}$ are given by the formula\n$$ \\partdiff{\\hat{F}}{x_i} = \\partdiff{F}{x_i}\n - \\phi(D({\\bf x})) \\partdiff{H}{x_i} - \\phi'(D({\\bf x})) \\partdiff{D}{x_i} H({\\bf x}) + 256 M|X|\\alpha \\partdiff{J}{x_i}. $$\nThe functions $F, H, D, J$ are infinitely differentiable, so the only possible issue is with $\\phi$.\nBy inspecting our construction of $\\phi$ (Lemma~\\ref{lemma:phi-construction}), we can see that it is piecewise \ninfinitely differentiable, and $\\phi'$ is continuous at the breakpoints. Therefore, it is also\nabsolutely continuous. This implies that $\\partdiff{\\hat{F}}{x_i}$ is absolutely continuous.\n\nThe function $\\hat{G}({\\bf x}) = F(\\bar{{\\bf x}}) + 256 M|X|\\alpha J({\\bf x})$ is infinitely differentiable,\nso its first partial derivatives are also absolutely continuous.\n\n\\item Assuming $\\partdiff{F}{x_i} \\geq 0$, we get $\\partdiff{\\tilde{F}}{x_i}\n \\geq - 64 M |X| \\alpha$ by Lemma~\\ref{lemma:F-bounds}. Using $\\partdiff{J}{x_i}\n \\geq |X|$, we get $\\partdiff{\\hat{F}}{x_i} = \\partdiff{\\tilde{F}}{x_i}\n + 256 M|X| \\alpha J({\\bf x}) \\geq 0$. The same holds for $\\partdiff{\\hat{G}}{x_i}$\nsince $\\partdiff{G}{x_i} \\geq 0$.\n\n\\item Assuming $\\mixdiff{F}{x_i}{x_j} \\leq 0$, we get $\\mixdiff{\\tilde{F}}{x_i}{x_j}\n \\leq 512 M |X| \\alpha$ by Lemma~\\ref{lemma:F-bounds}. Using $\\mixdiff{J}{x_i}{x_j}\n = -2$, we get $\\mixdiff{\\hat{F}}{x_i}{x_j} = \\mixdiff{\\tilde{F}}{x_i}{x_j}\n + 256 M |X| \\alpha \\mixdiff{J}{x_i}{x_j} \\leq 0$. The same holds for\n$\\mixdiff{\\hat{G}}{x_i}{x_j}$ since $\\mixdiff{G}{x_i}{x_j} \\leq 0$.\n\\end{enumerate}\n\\end{proof}\n\nThis concludes the proofs of our main hardness results.\n\n\n\n\\section{Algorithms using the multilinear relaxation}\n\\label{section:algorithms}\n\n\nHere we turn to our algorithmic results. First, we discuss\nthe problem of maximizing a submodular (but not necessarily\nmonotone) function subject to a matroid independence constraint.\n\n\n\\subsection{Matroid independence constraint}\n\\label{section:submod-independent}\n\nConsider the problem $\\max \\{ f(S): S \\in {\\cal I} \\}$, where ${\\cal I}$ is the collection\nof independent sets in a matroid ${\\cal M}$.\nWe design an algorithm based on the multilinear relaxation of the problem,\n$\\max \\{ F({\\bf x}): {\\bf x} \\in P({\\cal M}) \\}$.\nOur algorithm can be seen as \"continuous local search\" in the matroid polytope $P({\\cal M})$,\nconstrained in addition by the box $[0,t]^X$\nfor some fixed $t \\in [0,1]$. The intuition is that this forces\nour local search to use fractional solutions that are more fuzzy\nthan integral solutions and therefore less likely to get stuck in a local\noptimum. On the other hand, restraining the search space too much would not\ngive us much freedom in searching for a good fractional point.\nThis leads to a tradeoff and an optimal choice of $t \\in [0,1]$\nwhich we leave for later.\n\nThe matroid polytope is defined as\n$ P({\\cal M}) = \\mbox{conv} \\{ {\\bf 1}_I: I \\in {\\cal I} \\} $.\nor equivalently \\cite{E70} as\n$ P({\\cal M}) = \\{ {\\bf x} \\geq 0: \\forall S; \\sum_{i \\in S} x_i \\leq r_{\\cal M}(S) \\}$,\nwhere $r_{\\cal M}(S)$ is the rank function of ${\\cal M}$.\nWe define\n$$ P_t({\\cal M}) = P({\\cal M}) \\cap [0,t]^X = \\{ {\\bf x} \\in P({\\cal M}): \\forall i; x_i \\leq t \\}.$$\nWe consider the problem $\\max \\{F({\\bf x}): {\\bf x} \\in P_t({\\cal M})\\}$.\nWe remind the reader that\n$F({\\bf x}) = {\\bf E}[f(\\hat{{\\bf x}})]$ denotes the multilinear extension.\nOur algorithm works as follows.\n\n\\\n\n\\noindent{\\bf Fractional local search in $P_t({\\cal M})$} \\\\\n (given $t = \\frac{r}{q}$, $r \\leq q$ integer)\n\\begin{enumerate}\n\\item Start with ${\\bf x} := (0,0,\\ldots,0)$. Fix $\\delta = 1\/q$.\n\\item If there is $i,j \\in X$ and a direction ${\\bf v} \\in\n \\{{\\bf e}_j, -{\\bf e}_i, {\\bf e}_j - {\\bf e}_i \\}$ such that\n${\\bf x} + \\delta {\\bf v} \\in P_t({\\cal M})$ and $F({\\bf x} + \\delta {\\bf v}) > F({\\bf x})$,\nset ${\\bf x} := {\\bf x} + \\delta {\\bf v}$ and repeat.\n\\item If there is no such direction ${\\bf v}$, apply pipage rounding to ${\\bf x}$\nand return the resulting solution.\n\\end{enumerate}\n\n\\noindent{\\em Notes.}\nThe procedure as presented here would not run in polynomial\ntime. A modification which runs in polynomial time is that\nwe move to a new solution only if $F({\\bf x}+\\delta {\\bf v}) >\n F({\\bf x}) + \\frac{\\delta}{poly(n)} OPT$ (where we first get\na rough estimate of $OPT$ using previous methods).\nFor simplicity, we analyze the variant above\nand finally discuss why we can modify it without\nlosing too much in the approximation factor. We also defer\nthe question of how to estimate the value of $F({\\bf x})$\nto the end of this section.\n\nFor $t=1$, we have $\\delta=1$ and the procedure\nreduces to discrete local search. However, it is known that discrete\nlocal search alone does not give any approximation guarantee. With\nadditional modifications, an algorithm based on discrete local search\nachieves a $(\\frac14-o(1))$-approximation \\cite{LMNS09}.\n\nOur version of fractional local search avoids this issue and leads\ndirectly to a good fractional solution. Throughout the algorithm,\nwe maintain ${\\bf x}$ as a linear combination of $q$ independent sets\nsuch that no element appears in more than $r$ of them. A local step\ncorresponds to an add\/remove\/switch operation preserving this condition.\n\nFinally, we use pipage rounding to convert a fractional solution into\nan integral one. As we show in Lemma~\\ref{lemma:pipage} (in the Appendix),\na modification of the technique from \\cite{CCPV07}\ncan be used to find an integral solution without any loss in the objective\nfunction.\n\n\\begin{theorem}\n\\label{thm:submod-matroid-approx}\nThe fractional local search algorithm for any fixed $t \\in [0,\\frac12 (3-\\sqrt{5})]$\nreturns a solution of value at least $(t - \\frac12 t^2) OPT$,\nwhere $OPT = \\max \\{f(S): S \\in {\\cal I}\\}$.\n\\end{theorem}\n\nWe remark that for $t=\\frac12(3-\\sqrt{5})$, we would obtain a $\\frac14(-1+\\sqrt{5})\n \\simeq 0.309$-approximation, improving the factor of $\\frac14$ \\cite{LMNS09}.\nThis is not a rational value, but we can pick a rational $t$\narbitrarily close to $\\frac12 (3-\\sqrt{5})$.\nFor values $t > \\frac12 (3-\\sqrt{5})$, our analysis does not yield\na better approximation factor.\n\nFirst, we discuss properties of the point found by the fractional local search algorithm.\n\n\\begin{lemma}\n\\label{lemma:local-opt}\nThe outcome of the fractional local search algorithm $x$ is\na ``fractional local optimum'' in the following sense.\n(All the partial derivatives are evaluated at $x$.)\n\\begin{itemize}\n\\item For any $i$ such that ${\\bf x} - \\delta {\\bf e}_i \\in P_t({\\cal M})$,\n$\\partdiff{F}{x_i} \\geq 0.$\n\\item For any $j$ such that ${\\bf x} + \\delta {\\bf e}_j \\in P_t({\\cal M})$,\n$\\partdiff{F}{x_j} \\leq 0.$\n\\item For any $i,j$ such that ${\\bf x} + \\delta ({\\bf e}_j - {\\bf e}_i) \\in P_t({\\cal M})$,\n$\\partdiff{F}{x_j} - \\partdiff{F}{x_i} \\leq 0.$\n\\end{itemize}\n\\end{lemma}\n\n\\begin{proof}\nWe use the property (see \\cite{CCPV07}) that along any direction ${\\bf v} = \\pm {\\bf e}_i$\nor ${\\bf v} = {\\bf e}_i - {\\bf e}_j$, the function $F({\\bf x} + \\lambda {\\bf v})$ is a convex\nfunction of $\\lambda$. Also, observe that if it is possible to move from ${\\bf x}$\nin the direction of ${\\bf v}$ by any nonzero amount, then it is possible to move\nby $\\delta {\\bf v}$, because all coordinates of ${\\bf x}$ are integer multiples of $\\delta$\nand all the constraints also have coefficients which are integer multiples of $\\delta$.\nTherefore, if $\\frac{dF}{d\\lambda} > 0$ and it is possible\nto move in the direction of ${\\bf v}$, we would get\n$F({\\bf x} + \\delta {\\bf v}) > F({\\bf x})$ and the fractional local search would continue.\n\nIf ${\\bf v} = -{\\bf e}_i$ and it is possible to move along $-{\\bf e}_i$,\nwe get $\\frac{dF}{d\\lambda} = -\\partdiff{F}{x_i} \\leq 0$. Similarly,\nif ${\\bf v} = {\\bf e}_j$ and it is possible to move along ${\\bf e}_j$,\nwe get $\\frac{dF}{d\\lambda} = \\partdiff{F}{x_j} \\leq 0$. Finally,\nif ${\\bf v} = {\\bf e}_j - {\\bf e}_i$ and it is possible to move along ${\\bf e}_j - {\\bf e}_i$,\nwe get $\\frac{dF}{d\\lambda} = \\partdiff{F}{x_j} - \\partdiff{F}{x_i} \\leq 0$.\n\\end{proof}\n\nWe refer to the following exchange property for matroids\n(which follows easily from \\cite{Schrijver}, Corollary 39.12a; see also \\cite{LMNS09}).\n\n\\begin{lemma}\n\\label{lemma:exchange}\nIf $I, C \\in {\\cal I}$, then for any $j \\in C \\setminus I$,\nthere is $\\pi(j) \\subseteq I \\setminus C$, $|\\pi(j)| \\leq 1$,\nsuch that $I \\setminus \\pi(j) + j \\in {\\cal I}$.\nMoreover, the sets $\\pi(j)$ are disjoint\n(each $i \\in I \\setminus C$ appears at most once as $\\pi(j) = \\{i\\}$).\n\\end{lemma}\n\nUsing this, we prove a lemma about fractional local optima\nwhich generalizes Lemma 2.2 in \\cite{LMNS09}.\n\n\\begin{lemma}\n\\label{lemma:fractional-search}\nLet ${\\bf x}$ be the outcome of fractional local search over $P_t({\\cal M})$.\nLet $C \\in {\\cal I}$ be any independent set.\nLet $C' = \\{i \\in C: x_i < t\\}$. Then\n$$ 2 F({\\bf x}) \\geq F({\\bf x} \\vee {\\bf 1}_{C'}) + F({\\bf x} \\wedge {\\bf 1}_C).$$\n\\end{lemma}\n\nNote that for $t = 1$, the lemma reduces to\n$2 F({\\bf x}) \\geq F({\\bf x} \\vee {\\bf 1}_C) + F({\\bf x} \\wedge {\\bf 1}_C)$\n(similar to Lemma 2.2 in \\cite{LMNS09}).\nFor $t < 1$, however, it is necessary to replace $C$ by $C'$\nin the first expression, which becomes apparent in the proof.\nThe reason is that we do not have any information on $\\partdiff{F}{x_i}$\nfor coordinates where $x_i = t$.\n\n\\begin{proof}\nLet $C \\in {\\cal I}$ and assume ${\\bf x} \\in P_t({\\cal M})$ is a local optimum.\nSince ${\\bf x} \\in P({\\cal M})$, we can decompose it into a convex linear\ncombination of vertices of $P({\\cal M})$, ${\\bf x} = \\sum_{I \\in {\\cal I}} x_I {\\bf 1}_I$\nwhere $\\sum x_I = 1$..\nBy the smooth submodularity of $F({\\bf x})$ (see \\cite{Vondrak08}),\n\\begin{eqnarray*}\nF({\\bf x} \\vee {\\bf 1}_{C'}) - F({\\bf x}) \\leq \\sum_{j \\in C'} (1 - x_j) \\partdiff{F}{x_j} \n = \\sum_{j \\in C'} \\sum_{I: j \\notin I} x_I \\partdiff{F}{x_j}\n = \\sum_I x_I \\sum_{j \\in C' \\setminus I} \\partdiff{F}{x_j}.\n\\end{eqnarray*}\nAll partial derivatives here are evaluated at ${\\bf x}$.\nOn the other hand, also by submodularity,\n\\begin{eqnarray*}\n F({\\bf x}) - F({\\bf x} \\wedge {\\bf 1}_C) \\geq \\sum_{i \\notin C} x_i \\partdiff{F}{x_i}\n = \\sum_{i \\notin C} \\sum_{I: i \\in I} x_I \\partdiff{F}{x_i}\n = \\sum_I x_I \\sum_{i \\in I \\setminus C} \\partdiff{F}{x_i}.\n\\end{eqnarray*}\nTo prove the lemma, it remains to prove the following.\n\n\\\n\n\\noindent {\\bf Claim.} Whenever $x_I > 0$,\n$\\sum_{j \\in C' \\setminus I} \\partdiff{F}{x_j}\n \\leq \\sum_{i \\in I \\setminus C} \\partdiff{F}{x_i}$.\n\n\\\n\n\\noindent {\\em Proof:}\nFor any $I \\in {\\cal I}$, we can apply\nLemma~\\ref{lemma:exchange} to get a mapping $\\pi$ such that\n$I \\setminus \\pi(j) + j \\in {\\cal I}$ for any $j \\in C \\setminus I$.\nNow, consider $j \\in C' \\setminus I$, i.e. $j \\in C \\setminus I$ and $x_j < t$.\n\nIf $\\pi(j) = \\emptyset$, is possible to move from ${\\bf x}$ in the direction\nof ${\\bf e}_j$, because $I + j \\in {\\cal I}$ and hence we can replace $I$ by $I+j$\n(or at least we can do this for some nonzero fraction of its coefficient)\nin the linear combination.\nBecause $x_j < t$, we can move by a nonzero amount inside $P_t({\\cal M})$.\nBy Lemma~\\ref{lemma:local-opt}, $\\partdiff{F}{x_j} \\leq 0$.\n\nSimilarly, if $\\pi(j) = \\{i\\}$, it is possible to move in the direction\nof ${\\bf e}_j - {\\bf e}_i$, because $I$ can be replaced by $I \\setminus \\pi(j) + i$\nfor some nonzero fraction of its coefficient. By Lemma~\\ref{lemma:local-opt},\nin this case $\\partdiff{F}{x_j} - \\partdiff{F}{x_i} \\leq 0$.\n\nFinally, for any $i \\in I$ we have $x_i > 0$ and therefore we can decrease\n$x_i$ while staying inside $P_t({\\cal M})$. By Lemma~\\ref{lemma:local-opt},\nwe have $\\partdiff{F}{x_i} \\geq 0$ for all $i \\in I$.\nThis means\n$$ \\sum_{j \\in C' \\setminus I} \\partdiff{F}{x_j}\n \\leq \\sum_{j \\in C' \\setminus I: \\pi(j) = \\emptyset} \\partdiff{F}{x_j} +\n \\sum_{j \\in C' \\setminus I: \\pi(j) = \\{i\\}} \\partdiff{F}{x_i}\n \\leq \\sum_{i \\in I \\setminus C} \\partdiff{F}{x_i} $$\nusing the inequalities we derived above, and the fact that\neach $i \\in I \\setminus C$ appears at most once in $\\pi(j)$.\nThis proves the Claim, and hence the Lemma.\n\\end{proof}\n\nNow we are ready to prove Theorem~\\ref{thm:submod-matroid-approx}.\n\n\\begin{proof}\nLet ${\\bf x}$ be the outcome of the fractional local search over $P_t({\\cal M})$.\nDefine $A = \\{ i: x_i = t \\}$.\nLet $C$ be the optimum solution and $C' = C \\setminus A = \\{ i \\in C: x_i < t\\}$.\nBy Lemma~\\ref{lemma:fractional-search},\n$$ 2 F({\\bf x}) \\geq F({\\bf x} \\vee {\\bf 1}_{C'}) + F({\\bf x} \\wedge {\\bf 1}_C).$$\nFirst, let's analyze $F({\\bf x} \\wedge {\\bf 1}_C)$. We apply Lemma~\\ref{lemma:rnd-threshold} (in the Appendix),\nwhich states that $F({\\bf x} \\wedge {\\bf 1}_C) \\geq {\\bf E}[f(T({\\bf x} \\wedge {\\bf 1}_C))]$. Here,\n$T({\\bf x} \\wedge {\\bf 1}_C)$ is a random threshold set corresponding to the\nvector ${\\bf x} \\wedge {\\bf 1}_C$, i.e.\n$$ T({\\bf x} \\wedge {\\bf 1}_C) = \\{ i: ({\\bf x} \\wedge {\\bf 1}_C)_i > \\lambda \\}\n = \\{ i \\in C: x_i > \\lambda \\} = T({\\bf x}) \\cap C $$\nwhere $\\lambda \\in [0,1]$ is uniformly random.\nEquivalently,\n$$ F({\\bf x} \\wedge {\\bf 1}_C) \\geq {\\bf E}[f(T({\\bf x}) \\cap C)].$$\nDue to the definition of a threshold set, with probability $t$ we have\n$\\lambda < t$ and $T({\\bf x})$ contains $A = \\{ i: x_i = t \\} = C \\setminus C'$.\nThen, $f(T({\\bf x}) \\cap C) + f(C') \\geq f(C)$ by submodularity.\nWe conclude that\n\\begin{equation}\n\\label{eq:1}\nF({\\bf x} \\wedge {\\bf 1}_C) \\geq t (f(C) - f(C')).\n\\end{equation}\nNext, let's analyze $F({\\bf x} \\vee {\\bf 1}_{C'})$. We consider the ground set\npartitioned into $X = C \\cup \\bar{C}$, and we apply Lemma~\\ref{lemma:submod-split} (in the Appendix).\n(Here, $\\bar{C}$ denotes $X \\setminus C$, the complement of $C$ inside $X$.)\nWe get\n$$ F({\\bf x} \\vee {\\bf 1}_{C'}) \\geq {\\bf E}[f((T_1({\\bf x} \\vee {\\bf 1}_{C'}) \\cap C)\n \\cup (T_2({\\bf x} \\vee {\\bf 1}_{C'}) \\cap \\bar{C}))].$$\nThe random threshold sets look as follows: $T_1({\\bf x} \\vee {\\bf 1}_{C'}) \\cap C\n = (T_1({\\bf x}) \\cup C') \\cap C$ is equal to $C$ with probability $t$,\nand equal to $C'$ otherwise. $T_2({\\bf x} \\vee {\\bf 1}_{C'}) \\cap \\bar{C}\n = T_2({\\bf x}) \\cap \\bar{C}$ is empty with probability $1-t$. \n(We ignore the contribution when $T_2({\\bf x}) \\cap \\bar{C} \\neq \\emptyset$.)\nBecause $T_1$ and $T_2$ are independently sampled, we get\n$$ F({\\bf x} \\vee {\\bf 1}_{C'}) \\geq t(1-t) f(C) + (1-t)^2 f(C').$$\nProvided that $t \\in [0,\\frac12 (3 - \\sqrt{5})]$, we have $t \\leq (1-t)^2$.\nThen, we can write\n\\begin{equation}\n\\label{eq:1'}\nF({\\bf x} \\vee {\\bf 1}_{C'}) \\geq t(1-t) f(C) + t f(C').\n\\end{equation}\nCombining equations (\\ref{eq:1}) and (\\ref{eq:1'}), we get\n\\begin{eqnarray*}\n& & F({\\bf x} \\vee {\\bf 1}_{C'}) + F({\\bf x} \\wedge {\\bf 1}_C)\n \\geq t (f(C) - f(C')) + t(1-t) f(C) + t \\, f(C')\n = (2t - t^2) f(C).\n\\end{eqnarray*}\nTherefore,\n$$ F({\\bf x}) \\geq \\frac12 (F({\\bf x} \\vee {\\bf 1}_{C'}) + F({\\bf x} \\wedge {\\bf 1}_C))\n \\geq (t - \\frac12 t^2) f(C).$$\nFinally, we apply the pipage rounding technique which does not\nlose anything in terms of objective value (see Lemma~\\ref{lemma:pipage}).\n\\end{proof}\n\n\\paragraph{Technical remarks}\nIn each step of the algorithm, we need to estimate values of $F({\\bf x})$\nfor given ${\\bf x} \\in P_t({\\cal M})$. We accomplish this by using the expression\n$F({\\bf x}) = {\\bf E}[f(R({\\bf x}))]$ where $R({\\bf x})$ is a random set associated with ${\\bf x}$.\nBy standard bounds, if the values of $f(S)$ are in a range $[0,M]$,\nwe can achieve accuracy $M \/ poly(n)$ using a polynomial number of samples.\nWe use the fact that $OPT \\geq \\frac{1}{n} M$ (see\nLemma~\\ref{lemma:solution-value} in the Appendix) and therefore\nwe can achieve $OPT \/ poly(n)$ additive error in polynomial time.\n\nWe also relax the local step condition: we move to the next solution\nonly if $F({\\bf x} + \\delta {\\bf v}) > F({\\bf x}) + \\frac{\\delta}{poly(n)} OPT$\nfor a suitable polynomial in $n$. This way, we can only make a polynomial number of steps.\nWhen we terminate, the local optimality conditions (Lemma~\\ref{lemma:local-opt})\nare satisfied within an additive error of $OPT \/ poly(n)$,\nwhich yields a polynomially small error in the approximation bound.\n\n\n\\subsection{Matroid base constraint}\n\\label{section:submod-bases}\n\nLet us move on to the problem $\\max \\{f(S): S \\in {\\cal B}\\}$ where\n${\\cal B}$ are the bases of a matroid.\nFor a fixed $t \\in [0,1]$, let us consider an algorithm which can be seen\nas local search inside the base polytope $B({\\cal M})$, further constrained\nby the box $[0,t]^X$. The matroid base polytope is defined as\n$ B({\\cal M}) = \\mbox{conv} \\{ {\\bf 1}_B: B \\in {\\cal B} \\} $\nor equivalently \\cite{E70} as\n$ B({\\cal M}) = \\{ {\\bf x} \\geq 0: \\forall S \\subseteq X; \\sum_{i \\in S} x_i \\leq r_{\\cal M}(S),\n \\sum_{i \\in X} x_i = r_{\\cal M}(X) \\}, $\nwhere $r_{\\cal M}$ is the matroid rank function of ${\\cal M}$. Finally, we define\n$$ B_t({\\cal M}) = B({\\cal M}) \\cap [0,t]^X = \\{ {\\bf x} \\in B({\\cal M}): \\forall i \\in X; x_i \\leq t \\}.$$\nObserve that $B_t({\\cal M})$ is nonempty if and only if there is a convex linear\ncombination ${\\bf x} = \\sum_{B \\in {\\cal B}} \\xi_B {\\bf 1}_B$ such that $x_i \\in [0,t]$ for all $i$.\nThis is equivalent to saying that there is a linear combination (a fractional base packing)\n${\\bf x}' = \\sum_{B \\in {\\cal B}} \\xi'_B {\\bf 1}_B$ such that $x_i \\in [0,1]$\nand $\\sum \\xi'_B \\geq \\frac{1}{t}$, in other words the fractional base packing\nnumber is $\\nu \\geq \\frac{1}{t}$. Since the optimal fractional packing of bases\nin a matroid can be found efficiently (see Corollary 42.7a in \\cite{Schrijver}, Volume B),\nwe can find efficiently the minimum $t \\in [\\frac12,1]$ such that $B_t({\\cal M}) \\neq \\emptyset$.\nThen, our algorithm is the following.\n\n\\\n\n\\noindent{\\bf Fractional local search in $B_t({\\cal M})$} \\\\\n (given $t = \\frac{r}{q}$, $r \\leq q$ integer)\n\\begin{enumerate}\n\\item Let $\\delta = \\frac{1}{q}$.\nAssume that ${\\bf x} \\in B_t({\\cal M})$; adjust ${\\bf x}$ (using pipage rounding)\nso that each $x_i$ is an integer multiple of $\\delta$.\nIn the following, this property will be maintained.\n\\item If there is a direction ${\\bf v} = {\\bf e}_j - {\\bf e}_i$ such that\n${\\bf x} + \\delta {\\bf v} \\in B_t({\\cal M})$ and $F({\\bf x} + \\delta {\\bf v}) > F({\\bf x})$,\nthen set ${\\bf x} := {\\bf x} + \\delta {\\bf v}$ and repeat.\n\\item If there is no such direction ${\\bf v}$, apply pipage rounding to ${\\bf x}$\nand return the resulting solution.\n\\end{enumerate}\n\n\\noindent{\\em Notes.}\nWe remark that the starting point can be found as a convex linear combination\nof $q$ bases, $x = \\frac{1}{q} \\sum_{i=1}^{q} {\\bf 1}_{B_i}$, such that\nno element appears in more than $r$ of them, using matroid union techniques (see Theorem 42.9 in \\cite{Schrijver}).\nIn the algorithm, we maintain this representation.\nThe local search step corresponds to switching a pair of elements in one base,\nunder the condition that no element is used in more than $r$ bases at the same time.\nFor now, we ignore the issues of estimating $F({\\bf x})$ and stopping the local search\nwithin polynomial time. We discuss this at the end of this section.\n\nFinally, we use pipage rounding to convert the fractional solution ${\\bf x}$\ninto an integral one of value at least $F({\\bf x})$ (Lemma~\\ref{lemma:base-pipage} in the Appendix).\nNote that it is not necessarily true that any of the bases in\na convex linear combination ${\\bf x} = \\sum \\xi_B {\\bf 1}_{B}$ achieves the value $F({\\bf x})$.\n\n\\begin{theorem}\n\\label{thm:submod-bases-approx}\nIf there is a fractional packing of $\\nu \\in [1,2]$ bases in ${\\cal M}$,\nthen the fractional local search algorithm with $t = \\frac{1}{\\nu}$\nreturns a solution of value at least $\\frac12 (1-t) \\ OPT.$\n\\end{theorem}\n\nFor example, assume that ${\\cal M}$ contains two disjoint bases $B_1, B_2$\n(which is the case considered in \\cite{LMNS09}). Then,\nthe algorithm can be used with $t = \\frac{1}{2}$ and\nand we obtain a $(\\frac{1}{4}-o(1))$-approximation, improving the\n$(\\frac{1}{6}-o(1))$-approximation from \\cite{LMNS09}.\nIf there is a fractional packing of more than 2 bases,\nour analysis still gives only a $(\\frac{1}{4}-o(1))$-approximation.\nIf the dual matroid ${\\cal M}^*$ admits a better fractional packing of bases,\nwe can consider the problem $\\max \\{f(\\bar{S}): S \\in {\\cal B}^* \\}$ which\nis equivalent. For a uniform matroid, ${\\cal B} = \\{B: |B|=k\\}$,\nthe fractional base packing number is either at least $2$ or the same\nholds for the dual matroid, ${\\cal B}^* = \\{B: |B|=n-k\\}$ (as noted in \\cite{LMNS09}).\nTherefore, we get a $(\\frac{1}{4}-o(1))$-approximation for any uniform matroid.\nThe value $t=1$ can be used for any matroid,\nbut it does not yield any approximation guarantee.\n\n\n\n\n\\paragraph{Analysis of the algorithm}\nWe turn to the properties of fractional local optima.\nWe will prove that the point ${\\bf x}$ found by the fractional local search algorithm\nsatisfies the following conditions that allow us to compare $F({\\bf x})$ to the actual\noptimum.\n\n\\begin{lemma}\n\\label{lemma:local-opt2}\nThe outcome of the fractional local search algorithm ${\\bf x}$\nis a ``fractional local optimum'' in the following sense.\n\\begin{itemize}\n\\item For any $i,j$ such that ${\\bf x} + \\delta ({\\bf e}_j - {\\bf e}_i) \\in B_t({\\cal M})$,\n$$\\partdiff{F}{x_j} - \\partdiff{F}{x_i} \\leq 0.$$\n(The partial derivatives are evaluated at ${\\bf x}$.)\n\\end{itemize}\n\\end{lemma}\n\n\\begin{proof}\nSimilarly to Lemma~\\ref{lemma:local-opt},\nobserve that the coordinates of ${\\bf x}$ are always integer multiples of $\\delta$,\ntherefore if it is possible to move from ${\\bf x}$ in the direction of\n${\\bf v} = {\\bf e}_j - {\\bf e}_i$ by any nonzero amount, then it is possible to\nmove by $\\delta {\\bf v}$.\nWe use the property that for any direction ${\\bf v} = {\\bf e}_j - {\\bf e}_i$,\nthe function $F({\\bf x} + \\lambda {\\bf v})$ is a convex function of $\\lambda$ \\cite{CCPV07}.\nTherefore, if $\\frac{dF}{d\\lambda} > 0$ and it is possible\nto move in the direction of ${\\bf v}$, we would get\n$F({\\bf x} + \\delta {\\bf v}) > F({\\bf x})$ and the fractional local search would continue.\nFor ${\\bf v} = {\\bf e}_j - {\\bf e}_i$, we get\n$$ \\frac{dF}{d\\lambda} = \\partdiff{F}{x_j} - \\partdiff{F}{x_i} \\leq 0.$$\n\\end{proof}\n\nWe refer to the following exchange property for matroid bases\n(see \\cite{Schrijver}, Corollary 39.21a).\n\n\\begin{lemma}\n\\label{lemma:exchange2}\nFor any $B_1, B_2 \\in {\\cal B}$, there is a bijection $\\pi:B_1 \\setminus B_2\n \\rightarrow B_2 \\setminus B_1$ such that\n$\\forall i \\in B_1 \\setminus B_2$; $B_1-i+\\pi(i) \\in {\\cal M}$.\n\\end{lemma}\n\nUsing this, we prove a lemma about fractional local optima\nanalogous to Lemma 2.2 in \\cite{LMNS09}.\n\n\\begin{lemma}\n\\label{lemma:fractional-search2}\nLet ${\\bf x}$ be the outcome of fractional local search over $B_t({\\cal M})$.\nLet $C \\in {\\cal B}$ be any base. Then there is ${\\bf c} \\in [0,1]^X$ satisfying\n\\begin{itemize}\n\\item $c_i = t$, if $i \\in C$ and $x_i = t$\n\\item $c_i = 1$, if $i \\in C$ and $x_i < t$\n\\item $0 \\leq c_i \\leq x_i$, if $i \\notin C$\n\\end{itemize}\nsuch that\n$$ 2 F({\\bf x}) \\geq F({\\bf x} \\vee {\\bf c}) + F({\\bf x} \\wedge {\\bf c}).$$\n\\end{lemma}\n\nNote that for $t = 1$, we can set ${\\bf c} = {\\bf 1}_C$.\nHowever, in general we need this more complicated formulation.\nIntuitively, ${\\bf c}$ is obtained from ${\\bf x}$ by raising the variables $x_i, i \\in C$\nand decreasing $x_i$ for $i \\notin C$.\nHowever, we can only raise the variables $x_i, i \\in C$, where $x_i$ is below\nthe threshold $t$, otherwise we do not have any information about $\\partdiff{F}{x_i}$.\nAlso, we do not necessarily decrease all the variables outside of $C$ to zero.\n\n\\begin{proof}\nLet $C \\in {\\cal B}$ and assume ${\\bf x} \\in B_t({\\cal M})$ is a fractional local optimum.\nWe can decompose ${\\bf x}$ into a convex linear combination of vertices\nof $B({\\cal M})$, ${\\bf x} = \\sum \\xi_B {\\bf 1}_B$.\nBy Lemma~\\ref{lemma:exchange2}, for each base $B$ there is a bijection\n$\\pi_B:B \\setminus C \\rightarrow C \\setminus B$ such that\n$\\forall i \\in B \\setminus C$; $B - i + \\pi_B(i) \\in {\\cal B}$.\n\nWe define $C' = \\{i \\in C: x_i < t \\}$.\nThe reason we consider $C'$ is that\nif $x_i = t$, there is no room for an exchange step increasing $x_i$,\nand therefore Lemma~\\ref{lemma:local-opt2} does not give any information\nabout $\\partdiff{F}{x_i}$.\nWe construct the vector ${\\bf c}$ by starting from ${\\bf x}$, and for each $B$\nswapping the elements in $B \\setminus C$ for their image under $\\pi_B$,\nprovided it is in $C'$, until we raise the coordinates on $C'$ to $c_i=1$.\nFormally, we set $c_i = 1$ for $i \\in C'$, $c_i = t$ for $i \\in C \\setminus C'$,\nand for each $i \\notin C$, we define\n$$ c_i = x_i - \\sum_{B: i \\in B, \\pi_B(i) \\in C'} \\xi_B.$$\n\nIn the following, all partial derivatives are evaluated at ${\\bf x}$.\nBy the smooth submodularity of $F({\\bf x})$ (see \\cite{CCPV09}),\n\\begin{eqnarray}\n\\label{eq:2}\nF({\\bf x} \\vee {\\bf c}) - F({\\bf x}) & \\leq & \\sum_{j: c_j > x_j} (c_j - x_j) \\partdiff{F}{x_j}\n = \\sum_{j \\in C'} (1 - x_j) \\partdiff{F}{x_j}\n = \\sum_B \\sum_{j \\in C' \\setminus B} \\xi_B \\partdiff{F}{x_j}\n\\end{eqnarray}\nbecause $\\sum_{B: j \\notin B} \\xi_B = 1 - x_j$ for any $j$.\nOn the other hand, also by smooth submodularity,\n\\begin{eqnarray*}\nF({\\bf x}) - F({\\bf x} \\wedge {\\bf c}) & \\geq & \\sum_{i: c_i < x_i} (x_i - c_i) \\partdiff{F}{x_i}\n = \\sum_{i \\notin C} (x_i - c_i) \\partdiff{F}{x_i}\n = \\sum_{i \\notin C} \\sum_{B: i \\in B, \\pi_B(i) \\in C'} \\xi_B \\partdiff{F}{x_i}\n\\end{eqnarray*}\nusing our definition of $c_i$. In the last sum,\nfor any nonzero contribution, we have $\\xi_B > 0$, $i \\in B$ and $j = \\pi_B(i) \\in C'$,\ni.e. $x_j < t$. Therefore it is possible to move in the direction ${\\bf e}_j - {\\bf e}_i$\n(we can switch from $B$ to $B-i+j$).\nBy Lemma~\\ref{lemma:local-opt2},\n$$ \\partdiff{F}{x_j} - \\partdiff{F}{x_i} \\leq 0.$$\nTherefore, we get\n\\begin{eqnarray}\nF({\\bf x}) - F({\\bf x} \\wedge {\\bf c}) & \\geq &\n\\sum_{i \\notin C} \\sum_{B: i \\in B, j=\\pi_B(i) \\in C'} \\xi_B \\partdiff{F}{x_j}\n\\label{eq:3}\n = \\sum_B \\sum_{i \\in B \\setminus C: j=\\pi_B(i) \\in C'} \\xi_B \\partdiff{F}{x_j}.\n\\end{eqnarray}\nBy the bijective property of $\\pi_B$, this is equal to\n$\\sum_B \\sum_{j \\in C' \\setminus B} \\xi_B \\partdiff{F}{x_j}$.\nPutting (\\ref{eq:2}) and (\\ref{eq:3}) together, we get\n$F({\\bf x} \\vee {\\bf c}) - F({\\bf x}) \\leq F({\\bf x}) - F({\\bf x} \\wedge {\\bf c})$.\n\\end{proof}\n\nNow we are ready to prove Theorem~\\ref{thm:submod-bases-approx}.\n\n\\begin{proof}\nAssuming that $B_t({\\cal M}) \\neq \\emptyset$, we can find a starting point ${\\bf x}_0 \\in B_t({\\cal M})$.\nFrom this point, we reach a fractional local optimum ${\\bf x} \\in B_t({\\cal M})$\n(see Lemma~\\ref{lemma:local-opt2}).\nWe want to compare $F({\\bf x})$ to the actual optimum; assume that $OPT = f(C)$.\n\nAs before, we define $C' = \\{i \\in C: x_i < t\\}$.\nBy Lemma~\\ref{lemma:fractional-search2}, we know that\nthe fractional local optimum satisfies:\n\\begin{equation}\n\\label{eq:4}\n2 F({\\bf x}) \\geq F({\\bf x} \\vee {\\bf c}) + F({\\bf x} \\wedge {\\bf c})\n\\end{equation}\nfor some vector ${\\bf c}$ such that $c_i = t$ for all $i \\in C \\setminus C'$,\n$c_i = 1$ for $i \\in C'$ and $0 \\leq c_i \\leq x_i$ for $i \\notin C$.\n\nFirst, let's analyze $F({\\bf x} \\vee {\\bf c})$. We have\n\\begin{itemize}\n\\item $({\\bf x} \\vee {\\bf c})_i = 1$ for all $i \\in C'$.\n\\item $({\\bf x} \\vee {\\bf c})_i = t$ for all $i \\in C \\setminus C'$.\n\\item $({\\bf x} \\vee {\\bf c})_i \\leq t$ for all $i \\notin C$.\n\\end{itemize}\nWe apply Lemma~\\ref{lemma:submod-split} to the partition $X = C \\cup \\bar{C}$.\nWe get\n$$ F({\\bf x} \\vee {\\bf c}) \\geq {\\bf E}[f((T_1({\\bf x} \\vee {\\bf c}) \\cap C) \\cup (T_2({\\bf x} \\vee {\\bf c}) \\cap \\bar{C}))] $$\nwhere $T_1({\\bf x})$ and $T_2({\\bf x})$ are independent threshold sets.\nBased on the information above, $T_1({\\bf x} \\vee {\\bf c}) \\cap C = C$ with probability $t$\nand $T_1({\\bf x} \\vee {\\bf c}) \\cap C = C'$ otherwise. On the other hand,\n$T_2({\\bf x} \\vee {\\bf c}) \\cap \\bar{C} = \\emptyset$ with probability at least $1-t$.\nThese two events are independent. We conclude that on the right-hand side,\nwe get $f(C)$ with probability at least $t(1-t)$, or $f(C')$ with probability\nat least $(1-t)^2$:\n\\begin{equation}\n\\label{eq:5}\nF({\\bf x} \\vee {\\bf c}) \\geq t(1-t) f(C) + (1-t)^2 f(C').\n\\end{equation}\nTurning to $F({\\bf x} \\wedge {\\bf c})$, we see that\n\\begin{itemize}\n\\item $({\\bf x} \\wedge {\\bf c})_i = x_i$ for all $i \\in C'$.\n\\item $({\\bf x} \\wedge {\\bf c})_i = t$ for all $i \\in C \\setminus C'$.\n\\item $({\\bf x} \\wedge {\\bf c})_i \\leq t$ for all $i \\notin C$.\n\\end{itemize}\nWe apply Lemma~\\ref{lemma:submod-split} to $X = C \\cup \\bar{C}$.\n$$ F({\\bf x} \\wedge {\\bf c}) \\geq {\\bf E}[f((T_1({\\bf x} \\wedge {\\bf c}) \\cap C) \\cup (T_2({\\bf x} \\wedge {\\bf c}) \\cap \\bar{C}))].$$\nWith probability $t$, $T_1({\\bf x} \\wedge {\\bf c}) \\cap C$ contains $C \\setminus C'$\n(and maybe some elements of $C'$). In this case, $f(T_1({\\bf x} \\wedge {\\bf c}) \\cap C)\n \\geq f(C) - f(C')$ by submodularity.\nAlso, $T_2({\\bf x} \\wedge {\\bf c}) \\cap \\bar{C}$ is empty with probability at least $1-t$.\nAgain, these two events are independent.\nTherefore,\n$ F({\\bf x} \\wedge {\\bf c}) \\geq t(1-t) (f(C) - f(C')).$\nIf $f(C') > f(C)$, this bound is vacuous; otherwise, we can replace $t(1-t)$ by\n$(1-t)^2$, because $t \\geq 1\/2$. In any case,\n\\begin{equation}\n\\label{eq:6}\nF({\\bf x} \\wedge {\\bf c}) \\geq (1-t)^2 (f(C) - f(C')).\n\\end{equation}\nCombining (\\ref{eq:4}), (\\ref{eq:5}) and (\\ref{eq:6}),\n$$ F({\\bf x}) \\geq \\frac12 (F({\\bf x} \\vee {\\bf c}) + F({\\bf x} \\wedge {\\bf c}))\n \\geq \\frac12 (t(1-t) f(C) + (1-t)^2 f(C)) = \\frac12 (1-t) f(C).$$\n\\end{proof}\n\n\\paragraph{Technical remarks}\nAgain, we have to deal with the issues of estimating $F({\\bf x})$ and\nstopping the local search in polynomial time. We do this exactly\nas we did at the end of Section~\\ref{section:submod-independent}.\nOne issue to be careful about here is that if $f:2^X \\rightarrow [0,M]$,\nour estimates of $F({\\bf x})$ are within an additive error of $M \/ poly(n)$.\nIf the optimum value $OPT = \\max \\{f(S): S \\in {\\cal B}\\}$ is very small\ncompared to $M$, the error might be large compared to $OPT$ which would be\na problem. The optimum could in fact be very small in general.\nBut it holds that if ${\\cal M}$ contains no loops and co-loops\n(which can be eliminated easily), then $OPT \\geq \\frac{1}{n^2} M$\n(see Appendix~\\ref{app:base-value}). Then, our sampling errors are\non the order of $OPT \/ poly(n)$ which yields a $1\/poly(n)$ error\nin the approximation bound.\n\n\n\\section{Approximation for symmetric instances}\n\\label{section:symmetric}\n\nWe can achieve a better approximation assuming that the instance exhibits\na certain symmetry. This is the same kind of symmetry that we use in our hardness\nconstruction (Section~\\ref{section:hardness-proof}) and the hard instances exhibit\nthe same symmetry as well. It turns out that our approximation in this case\nmatches the hardness threshold up to lower order terms.\n\nSimilar to our hardness result, the symmetries that we consider\nhere are permutations of the ground set $X$, corresponding\nto permutations of coordinates in ${\\boldmath R}^X$. We start with some basic\nproperties which are helpful in analyzing symmetric instances.\n\n\\begin{lemma}\n\\label{lemma:grad-sym}\nAssume that $f:2^X \\rightarrow {\\boldmath R}$ is invariant with respect to\na group of permutations ${\\cal G}$ and $F({\\bf x}) = {\\bf E}[f(\\hat{{\\bf x}})]$.\nThen for any symmetrized vector $\\bar{{\\bf c}} = {\\bf E}_{\\sigma \\in {{\\cal G}}}[\\sigma({\\bf c})]$,\n$\\nabla F |_{\\bar{{\\bf c}}}$ is also symmetric w.r.t. ${\\cal G}$. I.e., for any $\\tau \\in {\\cal G}$,\n$$ \\tau(\\nabla F |_{{\\bf x}=\\bar{{\\bf c}}}) = \\nabla F |_{{\\bf x}=\\bar{{\\bf c}}}. $$\n\\end{lemma}\n\n\n\\begin{proof}\nSince $f(S)$ is invariant under ${\\cal G}$, so is $F({\\bf x})$,\ni.e. $F({\\bf x}) = F(\\tau({\\bf x}))$ for any $\\tau \\in {\\cal G}$.\nDifferentiating both sides at ${\\bf x}={\\bf c}$, we get by the chain rule:\n$$ \\partdiff{F}{x_i} \\Big|_{{\\bf x}={\\bf c}} =\n \\sum_j \\partdiff{F}{x_j} \\Big|_{{\\bf x}=\\tau({\\bf c})} \\partdiff{}{x_i} (\\tau(x))_j\n = \\sum_j \\partdiff{F}{x_j} \\Big|_{{\\bf x}=\\tau({\\bf c})} \\partdiff{x_{\\tau(j)}}{x_i}. $$\nHere, $\\partdiff{x_{\\tau(j)}}{x_i} = 1$ if $\\tau(j) = i$, and $0$ otherwise.\nTherefore,\n$$ \\partdiff{F}{x_i} \\Big|_{{\\bf x}={\\bf c}} =\n \\partdiff{F}{x_{\\tau^{-1}(i)}} \\Big|_{{\\bf x}=\\tau({\\bf c})}. $$\nNote that $\\tau(\\bar{{\\bf c}}) = {\\bf E}_{\\sigma \\in {\\cal G}}[\\tau(\\sigma({\\bf c}))]\n = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf c})] = \\bar{{\\bf c}}$\nsince the distribution of $\\tau \\circ \\sigma$ is equal to the distribution\nof $\\sigma$. Therefore,\n$$ \\partdiff{F}{x_i} \\Big|_{{\\bf x}=\\bar{{\\bf c}}} = \\partdiff{F}{x_{\\tau^{-1}(i)}}\n \\Big|_{{\\bf x}=\\bar{{\\bf c}}} $$\nfor any $\\tau \\in {\\cal G}$.\n\\end{proof}\n\n\nNext, we prove that the ``symmetric optimum''\n$\\max \\{F(\\bar{{\\bf x}}): {\\bf x} \\in P({\\cal F})\\}$ gives a solution which is\na local optimum for the original instance $\\max \\{F({\\bf x}): {\\bf x} \\in P({\\cal F})\\}$.\n(As we proved in Section~\\ref{section:hardness-proof}, in general\nwe cannot hope to find a better solution than the symmetric optimum.)\n\n\\begin{lemma}\n\\label{lemma:sym-local-opt}\nLet $f:2^X \\rightarrow {\\boldmath R}$ and ${\\cal F} \\subset 2^X$\nbe invariant with respect to a group of permutations $\\cal G$.\nLet $\\overline{OPT} = \\max \\{ F(\\bar{{\\bf x}}): {\\bf x} \\in P({\\cal F}) \\}$\nwhere $\\bar{{\\bf x}} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x})]$,\nand let ${\\bf x}_0$ be the symmetric point\nwhere $\\overline{OPT}$ is attained ($\\bar{{\\bf x}}_0 = {\\bf x}_0$).\nThen ${\\bf x}_0$ is a local optimum for the problem\n$\\max \\{F({\\bf x}): {\\bf x} \\in P({\\cal F}) \\}$, in the sense that\n$({\\bf x}-{\\bf x}_0) \\cdot \\nabla F |_{{\\bf x}_0} \\leq 0$ for any ${\\bf x} \\in P({\\cal F})$.\n\\end{lemma}\n\n\\begin{proof}\nAssume for the sake of contradiction that $({\\bf x}-{\\bf x}_0) \\cdot \\nabla F|_{{\\bf x}_0} > 0$\nfor some ${\\bf x} \\in P({\\cal F})$. We use the symmetric properties of $f$ and ${\\cal F}$\nto show that $(\\bar{{\\bf x}} - {\\bf x}_0) \\cdot \\nabla F|_{{\\bf x}_0} > 0$ as well.\nRecall that ${\\bf x}_0 = \\bar{{\\bf x}}_0$. We have\n$$ (\\bar{{\\bf x}} - {\\bf x}_0) \\cdot \\nabla F|_{{\\bf x}_0}\n = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x}-{\\bf x}_0) \\cdot \\nabla F|_{{\\bf x}_0}] \n = {\\bf E}_{\\sigma \\in {\\cal G}}[({\\bf x}-{\\bf x}_0) \\cdot \\sigma^{-1}(\\nabla F|_{{\\bf x}_0})]\n = ({\\bf x}-{\\bf x}_0) \\cdot \\nabla F|_{{\\bf x}_0} > 0 $$\nusing Lemma~\\ref{lemma:grad-sym}. Hence, there would be a direction\n$\\bar{{\\bf x}} - {\\bf x}_0$ along which an improvement can be obtained.\nBut then, consider a small $\\delta > 0$ such that\n${\\bf x}_1 = {\\bf x}_0 + \\delta (\\bar{{\\bf x}} - {\\bf x}_0) \\in P({\\cal F})$ and also\n$F({\\bf x}_1) > F({\\bf x}_0)$. The point ${\\bf x}_1$ is symmetric ($\\bar{{\\bf x}}_1 = {\\bf x}_1$)\nand hence it would contradict the assumption that $F({\\bf x}_0) = \\overline{OPT}$.\n\\end{proof}\n\n\n\\subsection{Submodular maximization over independent sets in a matroid}\n\nLet us derive an optimal approximation result for the problem\n$\\max \\{f(S): S \\in {\\cal I}\\}$ under the assumption that the instance is \"element-transitive\". \n\n\\begin{definition}\nFor a group ${\\cal G}$ of permutations on $X$, the orbit of an element $i \\in X$\nis the set $\\{ \\sigma(i): \\sigma \\in {\\cal G} \\}$.\n${\\cal G}$ is called element-transitive, if the orbit of any element is\nthe entire ground set $X$.\n\\end{definition}\n\nIn this case, we show that it is easy to achieve an optimal $(\\frac12-o(1))$-approximation\nfor a matroid independence constraint.\n\n\\begin{theorem}\n\\label{thm:sym-submod-independent}\nLet $\\max \\{f(S): S \\in {\\cal I}\\}$ be an instance symmetric with respect\nto an element-transitive group of permutations ${\\cal G}$. Let\n$\\overline{OPT} = \\max \\{F(\\bar{{\\bf x}}): {\\bf x} \\in P({\\cal M})\\}$\n where $\\bar{{\\bf x}} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x})]$.\nThen $\\overline{OPT} \\geq \\frac12 OPT$.\n\\end{theorem}\n\n\\begin{proof}\nLet $OPT = f(C)$.\nBy Lemma~\\ref{lemma:sym-local-opt}, $\\overline{OPT} = F({\\bf x}_0)$\nwhere ${\\bf x}_0$ is a local optimum for the problem $\\max \\{F({\\bf x}): {\\bf x} \\in P({\\cal M})$.\nThis means it is also a local optimum in the sense of Lemma~\\ref{lemma:local-opt},\nwith $t=1$.\nBy Lemma~\\ref{lemma:fractional-search},\n$$ 2 F({\\bf x}_0) \\geq F({\\bf x}_0 \\vee {\\bf 1}_C) + F({\\bf x}_0 \\wedge {\\bf 1}_C).$$\nAlso, ${\\bf x}_0 = \\bar{{\\bf x}}_0$. As we are dealing with an element-transitive\ngroup of symmetries, this means all the coordinates of ${\\bf x}_0$ are equal,\n${\\bf x}_0 = (\\xi,\\xi,\\ldots,\\xi)$. Therefore, ${\\bf x}_0 \\vee {\\bf 1}_C$ is equal\nto $1$ on $C$ and $\\xi$ outside of $C$. By Lemma~\\ref{lemma:rnd-threshold} (in the Appendix),\n$$ F({\\bf x}_0 \\vee {\\bf 1}_C) \\geq (1-\\xi) f(C).$$\nSimilarly, ${\\bf x}_0 \\wedge {\\bf 1}_C$ is equal to $\\xi$ on $C$ and $0$ outside of $C$.\nBy Lemma~\\ref{lemma:rnd-threshold},\n$$ F({\\bf x}_0 \\wedge {\\bf 1}_C) \\geq \\xi f(C).$$\nCombining the two bounds,\n$$ 2 F({\\bf x}_0) \\geq F({\\bf x}_0 \\vee {\\bf 1}_C) + F({\\bf x}_0 \\wedge {\\bf 1}_C) \n \\geq (1-\\xi) f(C) + \\xi f(C) = f(C) = OPT.$$\n\\end{proof}\n\nSince all symmetric solutions ${\\bf x} = (\\xi,\\xi,\\ldots,\\xi)$ form\na 1-parameter family, and $F(\\xi,\\xi,\\ldots,\\xi)$ is a concave function,\nwe can search for the best symmetric solution (within any desired accuracy)\nby binary search. By standard techniques, we get the following.\n\n\\begin{corollary}\nThere is a $(\\frac12 - o(1))$-approximation (\"brute force\" search\nover symmetric solutions) for the problem $\\max \\{f(S): S \\in {\\cal I}\\}$\nfor instances symmetric under an element-transitive group of permutations.\n\\end{corollary}\n\nThe hard instances for submodular maximization subject to a matroid\nindependence constraint correspond to refinements of the\nMax Cut instance for the graph $K_2$ (Section~\\ref{section:hardness-applications}).\nIt is easy to see that such instances are element-transitive,\nand it follows from Section~\\ref{section:hardness-proof} that a $(\\frac12+\\epsilon)$-approximation\nfor such instances would require exponentially many value queries.\nTherefore, our approximation for element-transitive instances is optimal.\n\n\n\\subsection{Submodular maximization over bases}\n\nLet us come back to the problem of submodular maximization over the bases\nof matroid. The property that $\\overline{OPT}$ is a local optimum\nwith respect to the original problem $\\max \\{F({\\bf x}): {\\bf x} \\in P({\\cal F})\\}$\nis very useful in arguing about the value of $\\overline{OPT}$.\nWe already have tools to deal with local optima from\nSection~\\ref{section:submod-bases}. Here we prove the following.\n\n\\begin{lemma}\n\\label{lemma:base-local-opt}\nLet $B({\\cal M})$ be the matroid base polytope of ${\\cal M}$ and\n${\\bf x}_0 \\in B({\\cal M})$ a local maximum for the submodular maximization problem\n$\\max \\{F({\\bf x}): {\\bf x} \\in B({\\cal M})\\}$, in the sense that $({\\bf x}-{\\bf x}_0) \\cdot \\nabla F|_{{\\bf x}_0} \\leq 0$\nfor any ${\\bf x} \\in B({\\cal M})$. Assume in addition that ${\\bf x}_0 \\in [s,t]^X$. Then\n$$ F({\\bf x}_0) \\geq \\frac12 (1-t+s) \\cdot OPT.$$\n\\end{lemma}\n\n\\begin{proof}\nLet $OPT = \\max \\{ f(B): B \\in {\\cal B} \\} = f(C)$.\nWe assume that ${\\bf x}_0 \\in B({\\cal M})$ is a local optimum with respect\nto any direction ${\\bf x}-{\\bf x}_0$, ${\\bf x} \\in B({\\cal M})$, so it is also a local optimum with respect\nto the fractional local search in the sense of Lemma~\\ref{lemma:fractional-search2},\nwith $t=1$. The lemma implies that\n$$ 2 F({\\bf x}_0) \\geq F({\\bf x} \\vee {\\bf 1}_C) + F({\\bf x} \\wedge {\\bf 1}_C).$$\nBy assumption, the coordinates of ${\\bf x} \\vee {\\bf 1}_C$ are equal to $1$ on $C$\nand at most $t$ outside of $C$. With probability $1-t$, a random threshold in $[0,1]$\nfalls between $t$ and $1$, and Lemma~\\ref{lemma:rnd-threshold} (in the Appendix) implies that\n$$ F({\\bf x} \\vee {\\bf 1}_C) \\geq (1-t) \\cdot f(C).$$\nSimilarly, the coordinates of ${\\bf x} \\wedge {\\bf 1}_C$ are $0$ outside of $C$,\nand at least $s$ on $C$. A random threshold falls between $0$ and $s$\nwith probability $s$, and Lemma~\\ref{lemma:rnd-threshold} implies that\n$$ F({\\bf x} \\wedge {\\bf 1}_C) \\geq s \\cdot f(C).$$\nPutting these inequalities together, we get\n$$2 F({\\bf x}_0) \\geq F({\\bf x} \\vee {\\bf 1}_C) + F({\\bf x} \\wedge {\\bf 1}_C) \\geq (1-t+s) \\cdot f(C).$$\n\\end{proof}\n\n\\\n\n\\noindent\n{\\bf Totally symmetric instances.}\nThe application we have in mind here is a special case of submodular\nmaximization over the bases of a matroid, which we call {\\em totally symmetric}.\n\n\\begin{definition}\n\\label{def:totally-symmetric}\nWe call an instance $\\max \\{f(S): S \\in {\\cal F} \\}$ totally symmetric with respect\nto a group of permutations ${\\cal G}$, if both $f(S)$ and ${\\cal F}$ are invariant\nunder ${\\cal G}$ and moreover, there is a point ${\\bf c} \\in P({\\cal F})$ such that\n${\\bf c} = \\bar{{\\bf x}} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf x})]$ for every ${\\bf x} \\in P({\\cal F})$.\nWe call ${\\bf c}$ the center of the instance.\n\\end{definition}\n\nNote that this is indeed stronger than just being invariant under ${\\cal G}$.\nFor example, an instance on a ground set $X = X_1 \\cup X_2$ could be symmetric\nwith respect to any permutation of $X_1$ and any permutation of $X_2$.\nFor any ${\\bf x} \\in P({\\cal F})$, the symmetric vector $\\bar{{\\bf x}}$ is constant\non $X_1$ and constant on $X_2$. However, in a totally symmetric instance,\nthere should be a unique symmetric point.\n\n\\paragraph{Bases of partition matroids}\nA canonical example of a totally symmetric instance is as follows.\nLet $X = X_1 \\cup X_2 \\cup \\ldots \\cup X_m$ and let integers $k_1,\\ldots,k_m$\nbe given. This defines a partition matroid ${\\cal M} = (X,{\\cal B})$, whose bases are\n$$ {\\cal B} = \\{B: \\forall j; |B \\cap X_j| = k_j \\}.$$\nThe associated matroid base polytope is\n$$ B({\\cal M}) = \\{{\\bf x} \\geq 0: \\forall j; \\sum_{i \\in X_j} x_i = k_j \\}.$$\nLet ${\\cal G}$ be a group of permutations such that the orbit of each element\n$i \\in X_j$ is the entire set $X_j$.\nThis implies that for any ${\\bf x} \\in B({\\cal M})$,\n$\\bar{{\\bf x}}$ is the same vector ${\\bf c}$, with coordinates $k_j \/ |X_j|$ on $X_j$.\nIf $f(S)$ is also invariant under ${\\cal G}$, we have a totally symmetric\ninstance $\\max \\{f(S): S \\in {\\cal B} \\}$. \n\n\\paragraph{Example: welfare maximization}\nTo present a more concrete example, consider $X_j = \\{ a_{j1}, \\ldots, a_{jk} \\}$\nfor each $j \\in [m]$, a set of bases ${\\cal B} = \\{B: \\forall j; |B \\cap X_j| = 1\\}$,\nand an objective function in the form $f(S) = \\sum_{i=1}^{k} v(\\{j: a_{ji} \\in S\\})$,\nwhere $v:2^{[m]} \\rightarrow {\\boldmath R}_+$ is a submodular function.\nThis is a totally symmetric instance, which captures the welfare maximization\nproblem for combinatorial auctions where each player has the same valuation function.\n(Including element $a_{ji}$ in the solution corresponds to allocating item $j$ to player $i$;\nsee \\cite{CCPV09} for more details.) We remark that here we consider possibly\nnonmonotone submodular functions, which is not common for combinatorial auctions;\nnevertheless the problem still makes sense.\n\n\\\n\nWe show that for such instances, the center point achieves an improved approximation.\n\n\\begin{theorem}\n\\label{thm:sym-solution}\nLet $\\max \\{f(S): S \\in {\\cal B}\\}$ be a totally symmetric instance.\nLet the fractional packing number of bases be $\\nu$ and the fractional\npacking number of dual bases $\\nu^*$. Then the center point ${\\bf c}$\nsatisfies\n$$ F({\\bf c}) \\geq \\left(1 - \\frac{1}{2\\nu} - \\frac{1}{2\\nu^*} \\right) OPT.$$\n\\end{theorem}\n\nRecall that in the general case, we get a $\\frac12 (1-1\/\\nu-o(1))$-approximation\n(Theorem~\\ref{thm:submod-bases-approx}). By passing to the dual matroid,\nwe can also obtain a $\\frac12 (1-1\/\\nu^*-o(1))$-approximation,\nso in general, we know how to achieve a $\\frac12 (1-1\/\\max\\{\\nu,\\nu^*\\}-o(1))$-approximation.\nFor totally symmetric instances where $\\nu=\\nu^*$, we improve this\nto the optimal factor of $1-1\/\\nu$.\n\n\\begin{proof}\nSince there is a unique center ${\\bf c} = \\bar{{\\bf x}}$ for any ${\\bf x} \\in B({\\cal M})$,\nthis means this is also the symmetric optimum $F({\\bf c}) = \\max \\{ F(\\bar{{\\bf x}}): {\\bf x} \\in B({\\cal M})\\}$.\nDue to Lemma~\\ref{lemma:sym-local-opt}, ${\\bf c}$ is a local optimum\nfor the problem $\\max \\{F({\\bf x}): {\\bf x} \\in B({\\cal M})\\}$.\n\nBecause the fractional packing number of bases is $\\nu$, we have $c_i \\leq 1\/\\nu$\nfor all $i$. Similarly, because the fractional packing number of dual bases\n(complements of bases) is $\\nu^*$,\nwe have $1-c_i \\leq 1\/\\nu^*$. This means that ${\\bf c} \\in [1-1\/\\nu^*,1\/\\nu]$.\nLemma~\\ref{lemma:base-local-opt} implies that\n$$ 2 F({\\bf c}) \\geq \\left( 1-\\frac{1}{\\nu} + 1-\\frac{1}{\\nu^*} \\right) OPT.$$\n\\end{proof}\n\n\\begin{corollary}\nLet $\\max \\{f(S): S \\in {\\cal B}\\}$ be an instance\non a partition matroid where every base takes at least\nan $\\alpha$-fraction of each part, at most a $(1-\\alpha)$-fraction of each part,\nand the submodular function $f(S)$ is invariant under a group ${\\cal G}$\nwhere the orbit of each $i \\in X_j$ is $X_j$.\nThen, the center point ${\\bf c} = {\\bf E}_{\\sigma \\in {\\cal G}}[\\sigma({\\bf 1}_B)]$\n(equal for any $B \\in {\\cal B}$) satisfies $F({\\bf c}) \\geq \\alpha \\cdot OPT$.\n\\end{corollary}\n\n\\begin{proof}\nIf the orbit of any element $i \\in X_j$ is the entire set $X_j$,\nit also means that $\\sigma(i)$ for a random $\\sigma \\in {\\cal G}$\nis uniformly distributed over $X_j$ (by the transitive property of ${\\cal G}$).\nTherefore, symmetrizing any fractional vector ${\\bf x} \\in B({\\cal M})$ gives\nthe same vector $\\bar{{\\bf x}} = {\\bf c}$, where $c_i = k_j \/ |X_j|$ for $i \\in X_j$.\nAlso, our assumptions mean that the fractional packing number of bases\nis $1\/(1-\\alpha)$, and the fractional packing number of dual bases is\nalso $1\/(1-\\alpha)$.\nDue to Lemma~\\ref{thm:sym-solution}, the center ${\\bf c}$ satisfies\n$F({\\bf c}) \\geq \\alpha \\cdot OPT$.\n\\end{proof}\n\nThe hard instances for submodular maximization over matroid bases\nthat we describe in Section~\\ref{section:hardness-applications} are exactly of this form\n(see the last paragraph of Section~\\ref{section:hardness-applications}, with $\\alpha=1\/k$).\nThere is a unique symmetric solution,\n$x = (\\alpha,\\alpha,\\ldots,\\alpha, 1-\\alpha,1-\\alpha,\\ldots,1-\\alpha)$.\nThe fractional base packing number for these matroids is\n$\\nu = 1\/(1-\\alpha)$ and Theorem~\\ref{thm:general-hardness} implies that\nany $(\\alpha+\\epsilon) = (1-1\/\\nu+\\epsilon)$-approximation for such matroids would\nrequire exponentially many value queries. Therefore, our approximation\nin this special case is optimal.\n\n\n\n\\section*{Acknowledgment}\nThe author would like to thank Jon Lee and Maxim Sviridenko for helpful discussions.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzgpia b/data_all_eng_slimpj/shuffled/split2/finalzzgpia new file mode 100644 index 0000000000000000000000000000000000000000..7707445bdbb3a3f228409fef9027b7c4906d2e03 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzgpia @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\n\n In \\cite{Barrett-Suli}, the authors established the existence of global-in-time weak solutions to the Navier--Stokes--Fokker--Planck equations arising in the kinetic theory of dilute polymer solutions and describing a large class of bead-spring chain models with finitely extensible nonlinear elastic (FENE) type spring potentials. For $\\O\\subset \\R^3$ a bounded domain, the solvent density $\\vr$ and the solvent velocity field $\\vu$ satisfy the following equations in the space-time cylinder $(0,T]\\times \\O$, $T>0$:\n\\be\\label{i1}\n\\d_t \\vr + \\Div_x (\\vr \\vu) = 0,\n\\ee\n\\be\\label{i2}\n \\d_t (\\vr\\vu)+ \\Div_x (\\vr \\vu \\otimes \\vu) +\\nabla_x p(\\vr) -\n\\Div_x \\SSS =\\Div_x \\TT +\\vr\\, \\ff,\n\\ee\nwhich we assume here to be supplemented with the no-slip boundary condition\n\\be\\label{i4}\n\\vu=\\mathbf{0}\\quad \\mbox{on}\\ (0,T]\\times \\partial\\Omega.\n\\ee\n\nIgnoring the effect of temperature changes, we consider a barotropic pressure law\n\\be\\label{pressure}\np = p(\\vr), \\ p(\\vr) \\approx \\vr^\\gamma \\ \\mbox{for large values of}\\ \\vr.\n\\ee\nThe {\\em Newtonian stress tensor} $\\SSS$ is defined by\n\\be\\label{i3}\n\\SSS = \\mu^S \\left( \\frac{\\nabla_x \\vu + \\nabla^{\\rm T}_x \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu) \\II \\right) + \\mu^B (\\Div_x \\vu) \\II,\n\\ee\nwith the shear and bulk viscosity coefficients $\\mu^S$ and $\\mu^B$ defined below (see \\eqref{mu-eta1}). In contrast with \\cite{Barrett-Suli}, where the shear and bulk viscosity coefficients are taken to be constant, $\\mu^S >0$ and $\\mu^B \\geq 0$, we consider here the case when they are functions of the polymer number density. Such an extension requires nontrivial modifications of the method used in \\cite{Barrett-Suli}, and represents the main contribution of the present paper.\n\nIn \\emph{a bead-spring chain model} consisting of $K+1$ beads coupled with $K$ elastic springs representing a polymer chain, the non-Newtonian elastic extra stress tensor $\\TT$ is defined by a version of the Kramers expression (cf. \\eqref{def-TT0} below), depending on the probability density function $\\psi$, which, in addition to $t$ and $x$, also depends on the conformation vector $(q_1^{\\rm T},\\dots q_K^{\\rm T})^{\\rm T} \\in \\R^{3K}$, with $q_i$ representing the $3$-component \\emph{conformation\/orientation vector} of the $i$th spring in the chain. Let $D:=D_1\\times \\cdots \\times D_K \\subset \\R^{3K}$ be the domain of admissible conformation vectors. Typically $D_i$ is the whole space $\\R^3$ or a bounded open ball centered at the origin $0$ in $\\R^3$, for each $i=1,\\dots,K$. When $K=1$, the model is referred to as the \\emph{dumbbell model}. Here we focus on finitely extensible nonlinear (or, briefly, FENE-type) bead-spring-chain models where $D_i=B(0,b_i^{\\frac12})$, a ball centered at the origin $0$ in $\\mathbb{R}^3$ and of radius $b_i^{\\frac12}$, with $b_i>0$ for each $i \\in \\{1,\\dots, K\\}$. The {\\em extra-stress tensor} $\\TT$ is defined by the formula:\n\\be\\label{def-TT0}\n\\TT (\\psi)(t,x) := \\TT_1 (\\psi) (t,x) -\\left(\\int_{D\\times D} \\gamma(q,q')\\, \\psi(t,x,q)\\,\\psi(t,x,q') \\, \\dq\\,\\dq'\\right)\\II,\n\\ee\nwhere, similarly to \\cite{Barrett-Suli}, the interaction kernel $\\g$ is assumed to be a positive constant $\\gamma(q,q')\\equiv \\de>0.$\nConsequently,\n\\be\\label{def-TT}\n\\TT (\\psi) := \\TT_1 (\\psi) -\\de \\left(\\int_{D} \\psi \\,\\dq\\right)^2\\II.\n\\ee\nThe first part, $\\TT_1(\\psi)$, of $\\TT(\\psi)$ is given by the {\\em Kramers expression}\n\\be\\label{def-TT1}\n\\t_1(\\psi) := k \\left[\\left(\\sum_{i=1}^K \\CC_i(\\psi)\\right)-(K+1) \\left(\\int_{D} \\psi \\ \\dq\\right)\\II\\right],\n\\ee\nwhere $k>0$ is the product of the Boltzmann constant and the absolute temperature and\n\\be\\label{def-Ci}\n\\CC_{i}(\\psi)(t,x) := \\int_D \\psi(t,x,q)\\, U_i'\\bigg(\\frac{|q_i|^2}{2}\\bigg)\\, q_i q_i^{\\rm T}\\,\\dq, \\quad i=1,\\dots,K.\n\\ee\n\nIn the expression \\eqref{def-Ci}, the smooth functions $U_i : [0,\\frac{b_i}{2}) \\to [0,\\infty)$, $i=1,\\dots,K,$ are the spring potentials satisfying\n$U_i(0)=0$, $\\lim_{s\\to \\frac{b_i}{2}-} U_i(s) =+\\infty$.\nWe introduce the {\\em partial Maxwellian} $M_i : D_i \\to [0,\\infty)$ by\n\\be\\label{def-Mi}\nM_i(q_i) := \\frac{1}{Z_i} {\\rm e}^{- U_i\\big(\\frac{|q_i|^2}{2}\\big)},\\quad \\mbox{where } Z_i:=\\int_{D_i} \\mathrm{e}^{- U_i\\big(\\frac{|p_i|^2}{2}\\big)}\\,{\\rm d}p_i.\n\\ee\nThe {\\em Maxwellian} $M : D \\to [0,\\infty)$ is then defined as the product of the $K$ partial Maxwellians: i.e., for any $q=(q_1^{\\rm T},\\dots, q_K^{\\rm T})^{\\rm T}$ in $D = D_1\\times\\cdots \\times D_K$, we have that\n\\[M(q) := \\prod_{i=1}^K M_i(q_i).\\]\nClearly,\n$\n\\int_D M(q)\\ \\dq=1.\n$\nBy direct calculations one verifies that, for any $i \\in \\{1,\\dots,K\\}$,\n\\be\\label{pt2-M}\nM(q)\\, \\nabla_{q_i} (M(q))^{-1} = - M(q)^{-1}\\, \\nabla_{q_i} M(q) = \\nabla_{q_i}\\left(U_i\\bigg(\\frac{|q_i|^2}{2}\\bigg)\\right) = U_i'\\bigg(\\frac{|q_i|^2}{2}\\bigg) q_i\n\\quad \\forall\\,q = (q_1^{\\rm T},\\dots,q_K^{\\rm T})^{\\rm T}\\in D.\n\\ee\n\nAs in \\cite{Barrett-Suli}, we shall suppose that, for any $i\\in\\{1,\\dots, K\\}$, there exist positive constants $c_{ij}, \\ j=1,\\dots,4$, and $\\th_i > 1$ such that\n\\ba\\label{pt1-Mi-Ui}\nc_{i1}\\left({\\rm dist}\\,(q_i,\\d D_i)\\right)^{\\th_i} \\leq M_i(q_i)\\leq c_{i2}\\left({\\rm dist}\\,(q_i,\\d D_i)\\right)^{\\th_i},\\quad\nc_{i3} \\leq \\left({\\rm dist}\\,(q_i,\\d D_i)\\right) U_i'\\left(\\frac{|q_i|^2}{2}\\right)\\leq c_{i4} \\quad \\forall\\,q_i\\in D_i.\n\\ea\nIt is then straightforward to deduce that\n\\be\\label{pt3-Mi-Ui}\n\\int_{D_i} \\left(1+ \\left( U_i\\bigg(\\frac{|q_i|^2}{2}\\bigg) \\right)^2+ \\left( U_i'\\bigg(\\frac{|q_i|^2}{2}\\bigg) \\right)^2 \\right) M_i(q_i)\\,\\dq_i <\\infty,\\quad i=1,\\dots, K.\n\\ee\n\n\n\nThe probability density function $\\psi$ satisfies the following {\\em Fokker--Planck equation} in $(0,T]\\times \\O\\times D$:\n\\ba\\label{eq-psi}\n\\d_t \\psi + \\Div_x (\\vu\\,\\psi) + \\sum_{i=1}^K \\Div_{q_i} \\left( (\\nabla_x \\vu)\\, q_i\\, \\psi \\right)= \\e \\Delta_x \\psi+\\frac{1}{4\\lambda} \\sum_{i=1}^K\\sum_{j=1}^K A_{ij}\\, \\Div_{q_i}\\left( M \\nabla_{q_j} \\left( \\frac{\\psi}{M} \\right)\\right).\n\\ea\n\nThe \\emph{centre-of-mass diffusion} term $\\e \\Delta_x \\psi$ is generally of the form\n$\n\\e \\Delta_x \\left( \\frac{\\psi}{\\zeta(\\vr)} \\right),\n$\nwhich involves the {\\em drag coefficient} $\\zeta(\\cdot)$ depending on the fluid density $\\vr$. Here we assume that $\\zeta$ is a constant function, which is, for simplicity, taken to be identically $1$. The constant parameter $\\e$ is the {\\em centre-of-mass diffusion coefficient}, which is strictly positive.\nThe positive parameter $\\l$ is called the \\emph{Deborah number}; it characterizes the elastic relaxation property of the fluid. The constant matrix $A=(A_{ij})_{1\\leq i,j\\leq K}$, called the \\textit{Rouse matrix}, is symmetric and positive definite. We denote by $A_0$ the smallest eigenvalue of $A$; clearly, $A_0>0$. We refer to Section 1 of Barrett and S\\\"uli \\cite{Barrett-Suli1} for a derivation of the Fokker--Planck equation \\eqref{eq-psi}.\n\nThe Fokker--Planck equation needs to be supplemented by suitable boundary conditions. For any $i=1,\\dots, K$, let $\\d \\overline D_i:= D_1\\times \\cdots \\times D_{i-1} \\times \\d D_i \\times D_{i+1} \\cdots \\times D_K$ and suppose that\n\\ba\\label{boundary-psi}\n\\left(\\frac{1}{4\\lambda} \\sum_{j=1}^K A_{ij}\\, \\Div_{q_i}\\left(M \\nabla_{q_j} \\left( \\frac{\\psi}{M} \\right)\\right) -(\\nabla_x \\vu) \\, q_i\\, \\psi \\right)\\cdot \\frac{q_i}{|q_i|}&=0 \\quad &&\\mbox{on}\\ (0,T]\\times \\O \\times \\d \\overline D_i,\\\\\n \\nabla_x\\psi \\cdot {\\bf n}&=0 \\quad &&\\mbox{on}\\ (0,T]\\times \\d\\O \\times D.\n\\ea\n\nFinally, we introduce the {\\em polymer number density} $\\eta$ defined by\n\\be\\label{def-eta}\n\\eta(t,x) := \\int_D \\psi(t,x,q)\\,\\dq, \\qquad (t,x)\\in [0,T]\\times \\O.\n\\ee\nBy (formally) integrating the partial differential equation \\eqref{eq-psi} over $D$ and using the boundary condition in $\\eqref{boundary-psi}_1$, and by integrating the boundary condition $\\eqref{boundary-psi}_2$ over $D$, we deduce the following partial differential equation and boundary condition for the function $\\eta$:\n\\be\\label{eq-eta}\n\\d_t \\eta + \\Div_x (\\vu \\,\\eta) = \\e \\Delta_x \\eta\\quad \\mbox{in}\\ (0,T] \\times \\O;\\qquad \\nabla_x \\eta \\cdot {\\bf n}=0 \\quad \\mbox{on}\\ (0,T]\\times \\d\\O.\n\\ee\nBy noting \\eqref{def-eta} we see that the expression for the extra-stress tensor in \\eqref{def-TT} and \\eqref{def-TT1} can also be expressed as follows:\n\\be\\label{def-TT-f}\n\\TT (\\psi) := k \\left(\\sum_{i=1}^K \\CC_i(\\psi)\\right)-\\left(k(K+1) \\eta + \\de\\, \\eta^2\\right) \\II.\n\\ee\nAs has already been pointed out above, our main objective is to consider a class of models of this form where the viscosity coefficients $\\mu^S = \\mu^S(\\eta)$ and $\\mu^B = \\mu^B(\\eta)$\ndepend on $\\eta$.\n\n\n\n\\subsection{Dissipative (finite-energy) weak solutions}\n\nWe adopt the following hypotheses on the initial data:\n\\ba\\label{ini-data}\n&\\vr(0,\\cdot) = \\vr_0(\\cdot) \\ \\mbox{with}\\ \\vr_0 \\geq 0 \\ {\\rm a.e.} \\ \\mbox{in} \\ \\O, \\quad \\vr_0 \\in L^\\gamma(\\O) \\ \\mbox{with}\\ \\gamma>\\textstyle{\\frac32};\\\\\n&\\vu(0,\\cdot) = \\vu_0(\\cdot) \\in L^r(\\O;\\R^3) \\ \\mbox{for some $r>1$}\\ \\mbox{such that}\\ \\vr_0|\\vu_0|^2 \\in L^1(\\O);\\\\\n&\\psi(0,\\cdot) = \\psi_0(\\cdot) \\ \\mbox{with}\\ \\psi_0 \\geq 0 \\ \\mbox{a.e. in} \\ \\O\\times D,\\quad \\psi_0 \\left(\\log \\frac{\\psi_0}{M}\\right) \\in L^1(\\O\\times D);\\\\\n&\\eta(0,\\cdot)=\\int_{D}\\psi_0 \\,\\dq =:\\eta_0 \\in L^2 (\\O).\n\\ea\nWe deduce from $\\eqref{ini-data}_1$ and $\\eqref{ini-data}_2$ by using H\\\"older's inequality that $(\\vr\\vu)(0,\\cdot) = \\vr_0 \\vu_0 =\\sqrt{\\vr_0}\\sqrt{\\vr_0} \\vu_0 \\in L^{\\frac{2\\g}{\\g+1}}(\\O;\\R^3)$.\n\n\\begin{definition}\\label{def-weaksl} We say that $(\\vr,\\vu,\\psi,\\eta)$ is a dissipative (finite-energy)\nweak solution in $(0,T]\\times \\O\\times D$ to the system of equations\n\\eqref{i1}--\\eqref{i3}, \\eqref{eq-psi}--\\eqref{def-TT-f}, supplemented by the initial data \\eqref{ini-data}, if:\n\\begin{itemize}\n\n\\item $\\vr \\geq 0 \\ {\\rm a.e.\\ in} \\ (0,T] \\times \\Omega,\\quad \\vr \\in C_w ([0,T]; L^\\gamma(\\Omega)),\\quad \\vu\\in L^{r}(0,T;W_0^{1,r}(\\Omega; \\R^3))\\quad \\mbox{for some $r>1$},$\n\\ba\\label{weak-est}\n&\\vr \\vu \\in C_w([0,T]; L^{\\frac{2 \\gamma}{ \\gamma + 1}}(\\Omega; \\R^3)),\\quad \\vr |\\vu|^2 \\in L^\\infty(0,T; L^{1}(\\Omega));\\\\\n&\\psi \\geq 0 \\ {\\rm a.e.\\ in} \\ (0,T] \\times \\Omega \\times D, \\quad \\psi \\in C_{w}([0,T];L^1(\\Omega\\times D)),\\\\\n& \\nabla_x \\psi\\in L^1((0,T)\\times \\O\\times D;\\R^3), \\quad M \\nabla_{q} \\left(\\frac{\\psi}{M}\\right)\\in L^1((0,T)\\times \\O\\times D;\\R^{3K}),\\\\\n&\\eta =\\int_D \\psi\\ \\dq \\ {\\rm a.e.\\ in} \\ (0,T] \\times \\Omega,\\quad \\eta \\in C_w ([0,T]; L^2(\\Omega)) \\cap L^2 (0,T; W^{1,2}(\\Omega)),\\\\\n&\\TT (\\psi) := k \\left(\\sum_{i=1}^K \\CC_i(\\psi)\\right)-\\left(k(K+1) \\eta + \\de\\, \\eta^2\\right) \\II \\quad {\\rm a.e.\\ in} \\ (0,T] \\times \\Omega,\\quad \\TT\\in L^1((0,T)\\times \\O;\\R^{3\\times 3}).\n\\ea\n\n\n\\item For any $t \\in (0,T]$ and any test function $\\phi \\in C^\\infty([0,T] \\times \\Ov{\\Omega})$, one has\n\n\\be\\label{weak-form1}\n\\int_0^t\\intO{\\big[ \\vr \\partial_t \\phi + \\vr \\vu \\cdot \\Grad \\phi \\big]} \\,\\dt' =\n\\intO{\\vr(t, \\cdot) \\phi (t, \\cdot) } - \\intO{ \\vr_{0} \\phi (0, \\cdot) },\n\\ee\n\\be\\label{weak-form2}\n\\int_0^t \\intO{ \\big[ \\eta \\partial_t \\phi + \\eta \\vu \\cdot \\Grad \\phi - \\e \\nabla_x\\eta \\cdot \\nabla_x \\phi \\big]} \\, \\dt' = \\intO{ \\eta(t, \\cdot) \\phi (t, \\cdot) } - \\intO{ \\eta_{0} \\phi (0, \\cdot) }.\n\\ee\n\n\\item For any $t \\in (0,T]$ and any test function $\\vvarphi \\in \\DC([0,T] \\times {\\Omega};\\R^3)$, one has\n\\ba\\label{weak-form3}\n&\\int_0^t \\intO{ \\big[ \\vr \\vu \\cdot \\partial_t \\vvarphi + (\\vr \\vu \\otimes \\vu) : \\Grad \\vvarphi + p(\\vr)\\, \\Div_x \\vvarphi - \\SSS : \\Grad \\vvarphi \\big] } \\, \\dt'\\\\\n&= \\int_0^t \\intO{ \\TT : \\nabla_x\\vvarphi - \\vr\\, \\ff \\cdot \\vvarphi} \\, \\dt' + \\intO{ \\vr \\vu (t, \\cdot) \\cdot \\vvarphi (t, \\cdot) } - \\intO{ \\vr_{0} \\vu_{0} \\cdot \\vvarphi(0, \\cdot) },\n\\ea\nwhere $\\SSS$ is defined by \\eqref{i3}.\n\n\\item For any $t \\in (0,T]$ and any test function $\\phi \\in \\DC([0,T] \\times \\Ov{\\Omega}\\times D)$, one has\n\\ba\\label{weak-form4}\n&\\int_0^t \\intO{\\int_D \\bigg[ \\psi \\partial_t \\phi + \\psi \\vu \\cdot \\Grad \\phi + \\sum_{i=1}^K (\\nabla_x \\vu)\\, q_i\\, \\psi \\cdot \\nabla_{q_i} \\phi - \\e\\nabla_x \\psi \\cdot \\nabla_x \\phi \\bigg]\\,\\dq}\\,\\dt'\\\\\n&=\\frac{1}{4\\lambda} \\sum_{i=1}^K\\sum_{j=1}^K A_{ij}\\int_0^t \\intO{ \\int_D M\\nabla_{q_j}\\left(\\frac{\\psi}{M}\\right)\\cdot \\nabla_{q_i}\\phi\\ \\dq} \\, \\dt' + \\intO{ \\int_D \\psi \\phi (t, \\cdot)- \\psi_{0} \\phi(0, \\cdot)\\,\\dq }.\n\\ea\n\n\\item The continuity equation holds in the sense of renormalized solutions:\n\\be\\label{weak-renormal}\n\\d_t b(\\vr) +\\Div_x (b(\\vr)\\vu) + (b'(\\vr)\\vr - b(\\vr))\\,\\Div_x \\vu =0 \\quad \\mbox{in} \\ \\mathcal{D}' ((0,T]\\times \\O),\n\\ee\nfor any\n\\be\\label{cond-b-renormal}\nb\\in C^1([0,\\infty)),\\quad |b'(s)s|+| b(s)| \\leq c<\\infty \\qquad \\forall\\, s\\in [0,\\infty).\n\\ee\n\\item Let $\\mathcal{F}(s) := s(\\log s -1)+1$ for $s>0$ and define $\\mathcal{F}(0):=\\lim_{s\\to 0+} \\mathcal{F}(s)=1$. For a.e. $t \\in (0,T]$, the\nfollowing \\emph{energy inequality} holds:\n\\ba\\label{energy}\n&\\int_{\\Omega} \\bigg[\n\\frac{1}{2} \\vr |\\vu|^2 + P(\\vr) +\\de\\, \\eta^2 + k \\int_D M \\mathcal{F} (\\widetilde \\psi)\\ \\dq\\bigg] (t, \\cdot) \\, \\dx\\\\\n&\\quad+ \\int_0^t \\intO{ \\mu^S \\bigg|\\frac{\\nabla \\vu + \\nabla^{\\rm T} \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu)\\, \\II\\, \\bigg|^2 +\\mu^B |\\Div_x \\vu|^2 } \\, \\dt' + 2\\, \\e\\, \\de \\int_0^t \\intO{ |\\nabla_x \\eta|^2} \\, \\dt'\\\\\n&\\quad+ \\e\\, k \\int_0^t \\intO{ \\int_D M\\bigg|\\nabla_x \\sqrt{\\widetilde \\psi}\\, \\bigg|^2 \\, \\dq} \\,\\dt'\n+ \\frac{k\\,A_0}{4\\lambda} \\int_0^t \\intO{ \\int_D M \\bigg|\\nabla_{q} \\sqrt{\\widetilde \\psi}\\,\\bigg|^2 \\, \\dq} \\, \\dt'\\\\\n&\\leq \\int_{\\Omega} \\bigg[ \\frac{1}{2} \\vr_{0} |\\vu_{0} |^2 + P(\\vr_0) + \\de\\, \\eta_0^2 + k \\int_D M \\mathcal{F} \\bigg(\\frac{\\psi_0}{M}\\bigg)\\, \\dq\\bigg] \\dx + \\int_0^t\\int_\\O \\vr \\,\\ff \\cdot \\vu \\,\\,\\dx\\, \\dt',\n\\ea\n where we have set $P(\\vr) := \\vr \\int_1^\\vr p(z)\/ z^2 \\, {\\rm d}z$ and $\\widetilde \\psi:=\\frac{\\psi}{M}$.\n\n\\end{itemize}\n\n\\end{definition}\n\n\n\n\\begin{remark} Definition \\ref{def-weaksl} is fairly standard. The energy\ninequality (\\ref{energy}) identifies an important class of weak solutions, usually termed \\emph{dissipative (finite-energy)}.\nWe note that, given a \\emph{smooth} solution, by tedious but rather straightforward calculations one can obtain the following a priori bound (see (1.22) in \\cite{Barrett-Suli}):\n\\ba\\label{energy0}\n&\\int_{\\Omega} \\left[\n\\frac{1}{2} \\vr |\\vu|^2 + P(\\vr) +\\de\\, \\eta^2 + k \\int_D M \\mathcal{F} (\\widetilde \\psi)\\, \\dq\\right] (t, \\cdot) \\, \\dx+ \\int_0^t \\intO{ \\SSS : \\Grad \\vu\\, } \\, \\dt' + 2\\, \\e\\, \\de \\int_0^t \\intO{ |\\nabla_x \\eta|^2}\\, \\dt'\\\\\n& \\quad + \\e\\, k \\int_0^t \\intO{ \\int_D M\\left|\\nabla_x \\sqrt{\\widetilde \\psi} \\,\\right|^2 \\dq} \\, \\dt' + \\frac{k}{4\\lambda} \\sum_{i=1}^K\\sum_{j=1}^K A_{ij}\\int_0^t \\intO{ \\int_D M \\nabla_{q_j} \\sqrt{\\widetilde \\psi} \\cdot \\nabla_{q_i} \\sqrt{\\widetilde \\psi} \\, \\dq} \\, \\dt' \\\\\n&\\leq \\int_{\\Omega} \\left[\n\\frac{1}{2} \\vr_{0} |\\vu_{0} |^2 + P(\\vr_0) + \\de\\, \\eta_0^2 + k \\int_D M \\mathcal{F} \\bigg(\\frac{\\psi_0}{M}\\bigg)\\, \\dq\\right]\\!\\dx + \\int_0^t\\int_\\O \\vr \\,\\ff \\cdot \\vu \\, \\dx\\,\\dt',\\nn\n\\ea\nwhich then implies (\\ref{energy}). Indeed, thanks to the form of the Newtonian stress tensor in \\eqref{i3}, direct calculations yield that\n$$\n\\SSS : \\Grad \\vu = \\mu^S \\left|\\frac{\\nabla \\vu + \\nabla^{\\rm T} \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu) \\II \\right|^2 +\\mu^B |\\Div_x \\vu|^2.\n$$\nHence, by the positive definiteness of the Rouse matrix $A=(A_{ij})_{1\\leq i,j\\leq K}$, and recalling that the smallest eigenvalue of $A$ is $A_0>0$, we deduce \\eqref{energy}.\n\\end{remark}\n\n\\subsection{Assumptions and main results}\n\nWe shall suppose that both $\\mu^S$ and $\\mu^B$ are $C^1$ functions of the polymer number density $\\eta$, and we adopt the following assumptions: there exist positive constants $c_j, \\ j=1,\\dots,5$, and an $\\o\\in\\R$ such that\n\\be\\label{mu-eta1}\nc_{1} (1+\\eta)^\\o \\leq \\mu^S(\\eta) \\leq c_2 (1+\\eta)^\\o, \\quad |(\\mu^S)'(\\eta)| \\leq c_3+c_4 (1+\\eta)^{\\o-1}, \\quad 0 \\leq \\mu^B(\\eta) \\leq c_5 (1+\\eta)^\\o\n\\quad \\forall\\,\\eta\\geq 0.\n\\ee\nIn addition, for the sake of simplicity, we shall assume that\n\\be\\label{press}\np(\\vr) = a \\vr^\\gamma, \\quad a > 0, \\ \\gamma > \\frac32.\n\\ee\n\nAs the complete proof of the existence of dissipative weak solutions in the special case of constant viscosity coefficients is already very long and technical (cf. \\cite{Barrett-Suli}), in the more general setting of polymer-number-dependent viscosity coefficients considered here we shall confine ourselves to establishing \\emph{weak sequential stability} of the family of dissipative weak solutions, whose existence we shall assume; we shall however indicate in Section \\ref{End} the main steps of a possible complete existence proof in the case of polymer-number-density-dependent viscosity coefficients.\n\nAccordingly, the main result of the paper reads as follows.\n\n\\begin{theorem}[Weak Sequential Stability]\\label{theorem} Let $\\{(\\vr_n,\\vu_n,\\psi_n,\\eta_n)\\}_{n\\in \\N}$ be a sequence of dissipative (finite-energy) weak solutions in the sense of Definition \\ref{def-weaksl} associated with the initial data $\\{(\\vr_{0,n},\\vu_{0,n},\\psi_{0,n},\\eta_{0,n})\\}_{n\\in\\N}$ satisfying:\n\\ba\\label{cond-ini-n}\n&\\vr_{0,n} \\geq 0 \\ \\mbox{a.e. in} \\ \\O, \\quad \\vr_{0,n} \\to \\vr_0 \\ \\mbox{strongly in}\\ L^\\g(\\O);\\\\\n&\\vu_{0,n} \\to \\vu_0 \\ \\mbox{in}\\ L^r(\\O;\\R^3) \\ \\mbox{for some $r>1$}\\ \\mbox{such that}\\ \\vr_{0,n}|\\vu_{0,n}|^2 \\to \\vr_0|\\vu_0|^2 \\ \\mbox{strongly in}\\ L^1(\\O);\\\\\n& \\psi_{0,n} \\geq 0 \\ \\mbox{a.e. in} \\ \\O\\times D,\\quad \\psi_{0,n} \\to \\psi_{0},\\quad \\psi_{0,n} \\left(\\log \\frac{\\psi_{0,n}}{M}\\right) \\to \\psi_0 \\left(\\log \\frac{\\psi_0}{M}\\right) \\ \\mbox{strongly in}\\ L^1(\\O\\times D);\\\\\n&\\eta_{0,n} = \\int_{D}\\psi_{0,n} \\, \\dq \\to \\eta_0 \\ \\mbox{strongly in} \\ L^2 (\\O).\n\\ea\n\nLet $\\vc{f} \\in L^\\infty((0,T) \\times \\Omega;\\R^3)$.\nSuppose that the exponent $\\g$ in \\eqref{press} and the parameter $\\o$ in \\eqref{mu-eta1} satisfy\n\\be\\label{mu-eta2}\n0\\leq \\o < \\frac{5}{3}, \\quad \\g>\\frac{3}{2}\n\\ee\nor\n\\be\\label{mu-eta3}\n -\\frac{4}{3}<\\o\\leq 0, \\quad \\g>\\frac{6}{4+3\\o}.\n\\ee\nThen, there exists a subsequence (not indicated) such that\n$$\n(\\vr_n,\\vu_n,\\psi_n,\\eta_n) \\to (\\vr,\\vu,\\psi,\\eta) \\quad \\mbox{as $n\\to \\infty$, in the sense of distributions (at least weakly in}\\ L^1),\n$$\nwhere the limit $(\\vr,\\vu,\\psi,\\eta)$ is a dissipative (finite-energy) weak solution in the sense of Definition \\ref{def-weaksl} associated with the initial data $(\\vr_0,\\vu_0,\\psi_0,\\eta_0)$.\n\n\n\\end{theorem}\n\nBefore embarking on the proof of Theorem \\ref{theorem} two remarks are in order.\n\n\\begin{remark}\\label{rem-thm1} The strong convergence assumptions in \\eqref{cond-ini-n} imply that\n\\ba\\label{cond-ini-n1}\n\\vr_{0,n} \\vu_{0,n} \\to \\vr_{0} \\vu_{0} \\quad \\mbox{strongly in}\\ L^{\\frac{2\\g}{1+\\g}}(\\O;\\R^3),\\qquad \\eta_0 = \\int_D \\psi_0\\, \\dq \\quad \\mbox{a.e. in}\\ \\O.\n\\ea\n\n\n\\end{remark}\n\n\\begin{remark}\\label{rem-thm2} It is important to note that we allow the viscosity coefficients $\\mu^B$ and $\\mu^S$ to decay to zero as the polymer number density tends to infinity; this is achieved at the expense of assuming a larger adiabatic exponent $\\g$; cf. \\eqref{mu-eta3}.\n\n\\end{remark}\n\n\n\nThe bulk of the rest of the paper is devoted to the proof of Theorem \\ref{theorem}. Comments on the possibility of carrying out a complete proof of the existence\nof dissipative (finite-energy) weak solutions are given in Section \\ref{End}. Throughout the rest of the paper, if there is no specification, $c$ will denote a positive constant depending only on the length $T$ of the time interval and the following quantity associated with the initial data:\n$$\n\\sup_{n\\in \\N}\\int_{\\Omega} \\bigg[\n\\frac{1}{2} \\vr_{0,n} |\\vu_{0,n} |^2 + P(\\vr_{0,n}) + \\de\\, \\eta_{0,n}^2 + k \\int_D M \\mathcal{F} \\bigg(\\frac{\\psi_{0,n}}{M}\\bigg)\\, \\dq\\bigg] \\dx.\n$$\nWe emphasize, however, that the value of $c$ may vary from line to line.\n\n\\section{Preliminaries}\n\nIn this section, we recall some concepts that will be used systematically throughout the rest of the paper, including Maxwellian-weighted Lebesgue and Sobolev spaces, embeddings of spaces of Banach-space-valued weakly-continuous functions, the Div-Curl lemma, and Riesz operators.\n\n\\subsection{Maxwellian-weighted spaces}\\label{sec:def-LMr-HM1}\nFor any $r \\in [1,\\infty)$, $L_M^r(D)$ denotes the Maxwellian-weighted Lebesgue space over $D$ with norm\n\\be\\label{def-LMD1}\n\\|u\\|_{L^r_M(D)}:=\\left(\\int_{ D} M |u(x)|^r\\, \\dq \\right)^{\\frac{1}{r}}.\n\\ee\nSimilarly, we define $L_M^r(\\O\\times D):=L^r(\\O;L^r_M(D))$ and the Maxwellian-weighted Sobolev spaces\n\\ba\\label{def-SMD2}\nH^1_M(D)&:=\\left\\{ u \\in L_{\\rm loc}^1( D) : \\|u\\|_{H^1_M( D)}^2:=\\int_{ D} M\\big( |u|^2 + |\\nabla_q u|^2 \\big)\\, \\dq<\\infty \\right\\}.\\\\\nH^1_M(\\O\\times D)&:=\\left\\{ u \\in L_{\\rm loc}^1(\\O\\times D) : \\|u\\|_{H^1_M(\\O\\times D)}^2:=\\int_{\\O\\times D} M\\big( |u|^2 + |\\nabla_x u|^2 +|\\nabla_q u|^2\\big)\\, \\dq \\ \\dx<\\infty \\right\\}.\\nonumber\n\\ea\nThe proof of the following lemma can be found in Appendix C and Appendix D in \\cite{Barrett-Suli2}.\n\\begin{lemma}\\label{lem:HM1-emb} The normed spaces $L^r_M (D)$, $L_M^r(\\O\\times D)$, $H^1_M(D)$ and $H^1_M(\\O \\times D)$ are Banach spaces.\nThe embedding $H^1(\\O;L_M^2(D)) \\hookrightarrow L^6(\\O;L_M^2(D))$ is continuous, and the embeddings\n$H^1_M(D) \\hookrightarrow L_M^2(D)$, $H^1_M(\\O\\times D) \\hookrightarrow L^2_M(\\O\\times D)$ are compact.\n\\end{lemma}\n\n\\subsection{On $C_w([0, T];X)$ type spaces}\n\nLet $X$ be a Banach space. We denote by $C_w([0, T];X)$ the set of all functions $u\\in L^\\infty(0, T;X)$ such that the mapping\n$t \\in [0, T] \\mapsto \\langle \\phi,u(t)\\rangle_{X}\\in \\R$ is continuous on $[0, T]$ for all $\\phi \\in X'$. Here and throughout the paper, we use $X'$ to denote the\ndual space of $X$, and $\\langle \\cdot,\\cdot \\rangle_{X}$ to denote the duality pairing between $X'$ and $X$.\n\nWhenever $X$ has a predual $E$, in the sense that $E'=X$, we denote by $C_{w*}([0, T];X)$ the set of all functions $u\\in L^\\infty(0, T;X)$ such that the mapping $t \\in [0, T] \\mapsto \\langle u(t),\\phi\\rangle_{E}\\in \\R$ is continuous on $[0, T]$ for all $\\phi \\in E$. We reproduce Lemma 3.1 from \\cite{Barrett-Suli}.\n\\begin{lemma}\\label{lem-Cw-Cw*}\nSuppose that $X$ and $Y$ are Banach spaces.\n\\begin{itemize}\n\\item[(i)] Assume that the space $X$ is reflexive and is continuously embedded in the space $Y$; then,\n$$ L^\\infty(0, T;X) \\cap C_w([0, T];Y) = C_w([0, T];X).$$\n\\item[(ii)] Assume that $X$ has a separable predual $E$ and $Y$ has a predual $F$ such that $F$ is continuously\nembedded in $E$; then,\n $$ L^\\infty(0, T;X) \\cap C_{w*}([0, T];Y) = C_{w*}([0, T];X).$$\n\\end{itemize}\n\n\\end{lemma}\n\nWe recall an Arzel\\`{a}--Ascoli type result in $C_w([0,T];L^s(\\O))$. We refer to Lemma 6.2 in \\cite{N-book} for the proof.\n\\begin{lemma}\\label{lem-Cw} Let $r, s \\in (1,\\infty)$ and let $\\O$ be a bounded Lipschitz domain in $\\R^d, \\ d\\geq 2$. Suppose $\\{g_n\\}_{n\\in \\N}$ is a sequence of functions in $C_w([0, T]; L^s(\\O))$ such that $\\{g_n\\}_{n\\in \\N}$ is bounded in $C([0, T]; W^{-1,r}(\\O)) \\cap \\ L^\\infty(0, T;L^s(\\O))$. Then, there exists a subsequence (not indicated) such that the following hold:\n\\begin{itemize}\n\\item[(i)] $g_n\\to g$ in $C_w([0, T]; L^s(\\O))$;\n\\item[(ii)] If, in addition, $r\\leq \\frac{d}{d-1}$, or $r>\\frac{d}{d-1}$ and $s>\\frac{d\\,r}{d+r}$, then $g_n\\to g$ strongly in $C([0, T]; W^{-1,r}(\\O))$.\n\\end{itemize}\n\\end{lemma}\n\n\\subsection{Div-Curl lemma}\n\nWe recall the celebrated Div-Curl lemma due to L. Tartar \\cite{L.Tartar}.\n\\begin{lemma}\\label{lem-div-curl2}\nLet $Q\\subset \\R^d$ be a domain and let $\\{({\\bf U}_n,{\\bf V}_n)\\}_{n\\in \\N}$ be a sequence of functions such that\n$$\n{\\bf U}_n \\to {\\bf U} \\ \\mbox{weakly in $L^p(Q;\\R^d)$}, \\quad {\\bf V}_n \\to {\\bf V} \\ \\mbox{weakly in $L^{q}(Q;\\R^d)$}, \\quad \\mbox{as $n\\to \\infty$},\n$$\nwhere\n$\n\\frac{1}{p}+\\frac 1q =\\frac{1}{r}<1.\n$\nSuppose in addition that, for some $s>0$,\n$$\n\\{\\Div \\,{\\bf U}_n\\}_{n\\in \\N} \\ \\mbox{is precompact in $W^{-1,s}(Q)$},\\quad \\{{\\rm curl} \\,{\\bf V}_n\\}_{n\\in \\N} \\ \\mbox{is precompact in $W^{-1,s}(Q;\\R^{d\\times d})$}.\n$$\nThen,\n$\n{\\bf U}_n \\cdot {\\bf V}_n \\to {\\bf U} \\cdot {\\bf V} \\ \\mbox{weakly in}\\ L^r(Q).\n$\n\n\\end{lemma}\n\n\\subsection{On Riesz type operators}\\label{sec:Riesz}\n\nThe Riesz operator $\\RR_{j}$, $1\\leq j \\leq d$, in $\\R^d$ is defined as a Fourier integral operator with symbol $\\frac{\\xi_j}{|\\xi|}$. That is, for any $u\\in \\mathcal{S}'(\\R^d)$, where $\\mathcal{S}'(\\R^d)$ denotes the space of tempered distributions on $\\R^d$, $\\RR_j$ is defined by\n$$\n\\RR_{j}[u] := \\mathfrak{F}^{-1}\\left[\\frac{\\xi_j}{|\\xi|} \\mathfrak{F}[u]\\right],\n$$\nwhere $\\mathfrak{F}$ is the Fourier transform and $\\mathfrak{F}^{-1}$ is the inverse Fourier transform. We then define, for any $u\\in \\mathcal{S}'(\\R^d)$,\n$$\n\\RR_{ij} [u]:=\\RR_i \\circ \\RR_j u= \\mathfrak{F}^{-1}\\left[\\frac{\\xi_i\\xi_j}{|\\xi|^2} \\mathfrak{F}[u]\\right],\\quad 1\\leq i,j\\leq d.\n$$\n\n\nWe define $\\A_j$ by\n$\n\\A_{j} [u] :=- \\mathfrak{F}^{-1}\\left[\\frac{i\\xi_j}{|\\xi|^2} \\mathfrak{F}[u]\\right].\n$\nSince the derivative $\\d_j$ and the Laplacian $\\Delta$ can be seen as Fourier integral operators with symbols $i\\xi_j$ and $-|\\xi|^{-2}$, respectively, we can write\n$$\n\\RR_{ij} = \\d_i\\d_j{\\Delta}^{-1},\\quad \\A_j = -\\d_{j} \\Delta^{-1}.\n$$\n\nLet the matrix-valued operator $\\RR$ and the vector-valued operator $\\A$ be defined by\n\\be\\label{def-Reize}\n\\RR=(\\RR_{ij})_{i,j=1}^d=(\\nabla\\otimes\\nabla) \\Delta^{-1},\\quad \\A:=(A_j)_{j=1}^d = - \\nabla \\Delta^{-1}.\n\\ee\nWe have\n$\n\\RR_{ij}= - \\d_i \\A_j,\\ \\sum_{j=1}^d \\RR_{jj}= -\\sum_{j}\\d_j\\A_j =\\II.\n$ By Theorem 1.55 and Theorem 1.57 in \\cite{N-book} the following result holds.\n\\begin{lemma}\\label{lem-Riesz1} For any $p \\in (1,\\infty)$ the operators $\\RR_j$ and $\\RR_{ij}$, $1\\leq i, j\\leq d$, are bounded from $L^p(\\R^d)$ to $L^p(\\R^d)$. That is,\nthere exists a positive constant $c=c(p,d)$ such that\n$$\n\\|\\RR_j [u] \\|_{L^p(\\R^d)}+ \\|\\RR_{ij} [u] \\|_{L^p(\\R^d)} \\leq c(p,d) \\|u \\|_{L^p(\\R^d)}\\qquad \\forall\\,u\\in L^p(\\R^d).\n$$\nMoreover, for any $u\\in L^p(\\R^d)$, $v\\in L^{p'}(\\R^d)$, $11$ satisfying\n$\n\\frac{1}{r}+\\frac 1p-\\frac 1d < \\frac 1s < \\min\\{1,\\frac{1}{r}+\\frac 1p\\},\n$\nthere exists a constant $c=c(r,p,s,d)$ such that:\n$$\n\\| \\RR_{ij}[w\\,v] - w \\RR_{ij}[v] \\|_{W^{\\b,s}(\\R^d)}\\leq c\\, \\|w\\|_{W^{1,r}(\\R^d)} \\|v\\|_{L^{p}(\\R^d)},\n$$\nfor any $i, j \\in \\{1, \\dots, d\\}$, where $\\b\\in (0,1)$ satisfies\n$\n\\frac{\\b}{d} =\\frac 1d + \\frac 1s - \\frac{1}{r} - \\frac 1p.\n$\n\n\n\\end{lemma}\n\nIn Lemma \\ref{lem-Riesz3}, we used the fractional-order Sobolev space $W^{\\b,s}(\\R^d)$. Let $\\O$ be the whole space $\\R^d$ or a bounded Lipschitz domain in $\\R^d$. For any $\\b \\in (0,1)$ and $s\\in [1,\\infty) $, we define\n\\be\\label{def-frac-Sob}\n W^{\\b,s}(\\O):=\\left\\{u\\in L^s(\\O) : \\| u \\|_{W^{\\b,s}(\\O)}:= \\| u \\|_{L^{s}(\\O)} + \\left(\\int_\\O\\int_\\O \\frac{|u(x)-u(y)|^s}{|x-y|^{d+\\b s}} \\ \\dx \\ {\\rm d}y\\right)^{\\frac{1}{s}} <\\infty\\right\\}.\n\\ee\nWe recall the following classical compact embedding theorem (see Theorem 7.1 in \\cite{NPV}).\n\\begin{lemma}\\label{lem-frac-sob} Let $\\O\\subset \\R^d$ be a bounded Lipschitz domain and suppose that $\\b \\in (0,1)$ and $s\\in [1,\\infty)$; then, the embedding\nof $W^{\\b,s}(\\O)$ into $L^s(\\O)$ is compact, i.e. $W^{\\b,s}(\\O)\\hookrightarrow \\hookrightarrow L^s(\\O)$.\n\n\\end{lemma}\n\n\n\n\\section{Uniform bounds}\\label{sec:estimate}\n\nLet $(\\vr,\\vu,\\psi,\\eta)$ be a dissipative (finite-energy) weak solution in the sense of Definition \\ref{def-weaksl} with initial data $(\\vr_0,\\vu_0,\\psi_0,\\eta_0)$ satisfying \\eqref{ini-data}. This section is devoted to establishing bounds on $(\\vr,\\vu,\\psi,\\eta)$ under the hypotheses \\eqref{mu-eta1}--\\eqref{mu-eta3}.\n\n\n\\subsection{Gronwall's inequality and uniform bounds}\n\nWe begin by noting that\n\\ba\\label{ini-est2}\n\\int_{\\Omega} \\int_D M \\mathcal{F} \\bigg(\\frac{\\psi_0}{M}\\bigg)\\,\\dq \\,\\dx =\\int_{\\Omega} \\int_D \\bigg( \\psi_0\\log\\bigg(\\frac{\\psi_0}{M}\\bigg)-\\psi_0 +M\\bigg)\\,\\dq \\,\\dx.\n\\ea\nAs $ s\\log s \\geq s-1$ for all $s\\geq 0$, it follows that $\\psi_0\\log\\big(\\frac{\\psi_0}{M}\\big)\\geq \\psi_0 - M$, which then implies that\n$$\n\\left|\\psi_0\\log\\bigg(\\frac{\\psi_0}{M}\\bigg)-\\psi_0 +M\\right|\\leq 2 \\left|\\psi_0\\log\\bigg(\\frac{\\psi_0}{M}\\bigg)\\right|.\n$$\nThus, by \\eqref{ini-data}, we have that\n\\be\\label{ini-estf}\n\\int_\\O \\bigg[\\frac{1}{2} \\vr_{0} |\\vu_{0} |^2 + P(\\vr_0) + \\de\\, \\eta_0^2 + k \\int_D M \\mathcal{F} \\bigg(\\frac{\\psi_0}{M}\\bigg)\\,\\dq\\bigg]\\dx \\leq c.\n\\ee\n\nSince $\\ff\\in L^\\infty((0,T)\\times \\O;\\R^3)$ and\n$$\n|\\vr \\,\\ff \\cdot \\vu |\\leq |\\ff| \\, |\\sqrt \\vr|\\,| \\sqrt\\vr \\vu| \\leq |\\ff|\\left( \\vr + \\vr|\\vu|^2 \\right) \\leq |\\ff|\\left(1 + \\vr^\\g + \\vr|\\vu|^2 \\right),\n$$\nby using Gronwall's inequality we deduce from \\eqref{energy} and \\eqref{ini-estf} that, for a.e. $t\\in (0,T)$,\n\\ba\\label{energy-f}\n&\\int_{\\Omega} \\bigg[ \\frac{1}{2} \\vr |\\vu|^2 + P(\\vr) +\\de\\, \\eta^2 + k \\int_D M \\mathcal{F} (\\widetilde \\psi)\\, \\dq\\bigg] (t, \\cdot) \\,\\dx\\\\\n&\\quad+ \\int_0^t \\intO{ \\mu^S(\\eta) \\left|\\frac{\\nabla \\vu + \\nabla^{\\rm T} \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu) \\II \\right|^2 +\\mu^B(\\eta) |\\Div_x \\vu|^2 } \\, \\dt' + 2\\, \\e\\, \\de \\int_0^t \\intO{ |\\nabla_x \\eta|^2} \\, \\dt' \\\\\n&\\quad+\\e\\, k \\int_0^t \\intO{ \\int_D M\\left|\\nabla_x \\sqrt{\\widetilde \\psi}\\, \\right|^2\\dq} \\, \\dt+\\frac{k\\,A_0}{4\\lambda} \\int_0^t \\intO{ \\int_D M \\left|\\nabla_{q} \\sqrt{\\widetilde \\psi}\\,\\right|^2\\dq} \\, \\dt' \\\\\n&\\leq c(1+T)\\, {\\rm e}^{ct}.\n\\ea\nIn the rest of this section, we shall establish additional bounds on the unknowns, one by one, by using \\eqref{energy-f}.\n\n\\subsection{Bounds on the fluid density and the polymer number density}\n\nFrom \\eqref{energy-f}, we have\n\\be\\label{est-r-eta1}\n\\vr\\in L^\\infty(0,T;L^\\gamma(\\O)), \\quad \\eta\\in L^\\infty(0,T;L^2(\\O))\\cap L^2(0,T;W^{1,2}(\\O)), \\quad \\vr |\\vu|^2 \\in L^{\\infty}(0,T;L^{1}(\\O)).\n\\ee\nBy Sobolev embedding and interpolation it follows that\n\\be\\label{est-r-eta3}\n \\eta\\in L^\\infty(0,T;L^2(\\O)) \\cap L^2(0,T;L^6(\\O)) \\hookrightarrow L^{a}(0,T;L^{\\frac{6a}{3a-4}}(\\O)),\\quad 2\\leq a\\leq \\infty.\n \\ee\nThe bounds in \\eqref{est-r-eta1} then imply that\n\\be\\label{est-ru2}\n\\vr \\vu = \\sqrt{\\vr} \\sqrt{\\vr}\\vu \\in L^{\\infty}(0,T;L^{\\frac{2\\g}{\\g+1}}(\\O;\\R^3)).\n\\ee\n\n\\subsection{Bounds on the fluid velocity field}\\label{sec:est-u}\n\nBy \\eqref{energy-f} we deduce that\n\\be\\label{est-u1}\ng:= \\sqrt{\\mu^S(\\eta)} \\left|\\frac{\\nabla \\vu + \\nabla^{\\rm T} \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu) \\II \\right| \\in L^2(0,T;L^2(\\O)).\n\\ee\nWe recall that $\\mu^S(\\eta)$ fulfills the hypotheses stated in \\eqref{mu-eta1} with $\\o$ satisfying \\eqref{mu-eta2} or \\eqref{mu-eta3}. We begin by considering the case when $\\o\\geq 0$; \\eqref{mu-eta1} then implies that $\\mu^{-1}$ is uniformly bounded. Thus,\n\\be\\label{est-u2}\n \\left|\\frac{\\nabla \\vu + \\nabla^{\\rm T} \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu) \\II \\right| \\in L^2(0,T;L^2(\\O));\n\\ee\nwhence\nKorn's inequality (see \\cite{Dain}) with the no-slip boundary condition on $\\vu$ implies that\n\\be\\label{est-u4}\n|\\nabla_x \\vu| \\in L^2(0,T;L^2(\\O)), \\qquad \\mbox{and hence}\\qquad \\vu \\in L^2(0,T;W^{1,2}_0(\\O;\\R^3)).\n\\ee\nOn the other hand, when $\\o\\leq 0$ satisfies the constraint \\eqref{mu-eta3}, that is $-4\/3<\\o\\leq 0$, then, by \\eqref{mu-eta1} and \\eqref{est-r-eta3}, we have that\n\\be\\label{est-mu1}\n\\mu(\\eta)^{-1}\\in L^\\infty(0,T;L^{\\frac{2}{|\\o|}}(\\O)) \\cap L^{\\frac{2}{|\\o|}}(0,T;L^{\\frac{6}{|\\o|}}(\\O)) \\cap L^{\\frac{10}{3|\\o|}}((0,T)\\times \\O).\n\\ee\nThus, from \\eqref{est-u1}, we deduce that\n\\be\\label{est-u11}\n \\left|\\frac{\\nabla \\vu + \\nabla^{\\rm T} \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu) \\II \\right| \\in L^2(0,T;L^{\\frac{4}{2+|\\o|}}(\\O)) \\cap L^{\\frac{4}{2+|\\o|}}(0,T;L^{\\frac{12}{6+|\\o|}}(\\O))\\cap L^{\\frac{20}{10+3|\\o|}}((0,T)\\times \\O).\n\\ee\nHence, taking advantage of the no-slip boundary condition, we may use\nKorn's inequality to obtain\n\\be\\label{est-u13}\n\\vu \\in L^2(0,T;W_0^{1,\\frac{4}{2+|\\o|}}(\\O;\\R^3)) \\cap L^{\\frac{4}{2+|\\o|}}(0,T;W_0^{1,\\frac{12}{6+|\\o|}}(\\O;\\R^3))\\cap L^{\\frac{20}{10+3|\\o|}}(0,T;W_0^{1,\\frac{20}{10+3|\\o|}}(\\O;\\R^3)).\n\\ee\n\n\\subsection{Bounds on the Newtonian stress tensor}\n\nConsider first the case when \\eqref{mu-eta2} holds with $\\o\\geq 0$. Let us write\n$$\n\\mu^S(\\eta) \\left|\\frac{\\nabla \\vu + \\nabla^{\\rm T} \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu) \\II \\right| = g \\sqrt{\\mu^S(\\eta)},\n$$\nwhere $g\\in L^2((0,T)\\times \\O)$ is defined by \\eqref{est-u1}. Thanks to \\eqref{mu-eta1}, \\eqref{est-r-eta3}, \\eqref{est-u4}, and by a similar argument as in the derivation of \\eqref{est-u13}, we obtain\n\\be\\label{est-SS1}\n\\SSS\\in L^2(0,T;L^{\\frac{4}{2+|\\o|}}(\\O;\\R^{3\\times 3})) \\cap L^{\\frac{4}{2+|\\o|}}(0,T;L^{\\frac{12}{6+|\\o|}}(\\O;\\R^{3\\times 3}))\\cap L^{\\frac{20}{10+3|\\o|}}((0,T)\\times \\O;\\R^{3\\times 3}).\n\\ee\nOn the other hand, when $\\o\\leq 0$, by observing that $\\sqrt{\\mu^S} + \\sqrt{\\mu^B}\\leq c<\\infty$ we deduce that\n\\be\\label{est-SS2}\n\\SSS\\in L^2((0,T)\\times\\O;\\R^{3\\times 3}).\n\\ee\n\n\n\\subsection{Bounds on the probability density function}\\label{sec:est-psi}\n\nIt follows from \\eqref{energy-f} that\n\\be\\label{est-psi1}\n\\mathcal{F}(\\tpsi) \\in L^\\infty(0,T;L^1_M(\\O\\times D)), \\quad \\sqrt{\\tpsi} \\in L^2(0,T;H_M^1(\\O\\times D)),\n\\ee\nwhere we recall that $\\tpsi := \\psi\/M$. Consequently, by Sobolev embedding and thanks to Lemma \\ref{lem:HM1-emb} we then have that $\\sqrt{\\tpsi} \\in L^2(0,T;L^6(\\O;L^2_M(D)))$.\n\nIf $s>{\\rm e}^2$, then $s\\log s > 2s >2(s-1)$, which then implies that $\\mathcal{F}(s)=s\\log s -(s-1)>\\frac{1}{2}{s\\log s}$ for $s>{\\rm e}^2$. Thus,\n\\ba\\label{est-psi3}\n\\|\\tpsi\\log\\tpsi\\|_{L^1_M(\\O\\times D))}&=\\int_{\\O\\times D} |\\tpsi\\log\\tpsi|\\,M\\,\\dq\\,\\dx = \\int_{0 \\leq \\tpsi \\leq {\\rm e}^2} |\\tpsi\\log\\tpsi|\\,M\\,\\dq\\,\\dx +\\int_{\\tpsi > e^2} |\\tpsi\\log\\tpsi|\\,M\\,\\dq\\,\\dx \\\\\n& \\leq 2 {\\rm e}^2\\int_{\\O\\times D}\\,M\\,\\dq\\,\\dx + 2 \\int_{\\tpsi > {\\rm e}^2} |\\mathcal{F}(\\tpsi)|\\, M\\,\\dq\\,\\dx \\leq c,\n\\ea\nwhere we have used \\eqref{pt1-Mi-Ui} and \\eqref{est-psi1}. This implies that\n\\be\\label{est-psi5}\n\\tpsi\\log\\tpsi \\in L^\\infty(0,T;L^1_M(\\O\\times D)),\\quad \\psi\\log\\psi \\in L^\\infty(0,T;L^1(\\O\\times D)).\n\\ee\n\nClearly,\n\\be\\label{nabla-psi}\n\\nabla_{x}\\tpsi = 2 \\sqrt{\\tpsi} \\ \\nabla_{x}\\sqrt{\\tpsi},\\quad \\nabla_{q}\\tpsi = 2 \\sqrt{\\tpsi} \\ \\nabla_{q}\\sqrt{\\tpsi}.\n\\ee\nBy \\eqref{nabla-psi} and H\\\"older's inequality we therefore have that\n\\ba\\label{est-psi8}\n\\|\\nabla_{q,x} \\tpsi\\|_{L^1_M(D;\\R^{3(K+1)})} & = \\int_{D} | M\\,\\nabla_{q,x}\\tpsi | \\, \\dq = 2\\int_{D} M \\sqrt{\\tpsi} \\, \\left|\\nabla_{q,x}\\sqrt{\\tpsi}\\,\\right|\\dq \\\\\n&\\leq 2 \\left(\\int_{D} M \\left|\\nabla_{q,x}\\sqrt{\\tpsi}\\,\\right|^2 \\, \\dq \\right)^{\\frac 12}\\left(\\int_{D} M \\tpsi \\, \\dq \\right)^{\\frac 12}=2 \\left\\|\\nabla_{q,x}\\sqrt{\\tpsi}\\,\\right\\|_{L^2_M(D;\\R^{3(K+1)})} \\eta^{\\frac 12}.\n\\ea\nThus, we benefit from the estimates in \\eqref{est-r-eta3} for $\\eta$ and deduce that\n\\ba\\label{est-psi9}\n&\\nabla_{q,x} \\tpsi \\in L^2(0,T;L^{\\frac 43}(\\O;L^1_M(D;\\R^{3(K+1)})))\\cap L^{\\frac 43}(0,T;L^\\frac{12}{7}(\\O;L^1_M(D;\\R^{3(K+1)}))).\n\\ea\nBy observing that $\\|\\nabla_x \\tpsi\\|_{L^1_M(D;\\R^{3})}=\\|\\nabla_x\\psi\\|_{L^1(D;\\R^{3})}$ we obtain\n\\ba\\label{est-psi10}\n\\nabla_x \\psi \\in L^2(0,T;L^{\\frac 43}(\\O;L^1(D;\\R^{3})))\\cap L^{\\frac 43}(0,T;L^\\frac{12}{7}(\\O;L^1(D;\\R^{3}))).\n\\ea\n\n\n\n\\subsection{Bounds on the extra-stress tensor}\n\nWe recall that the extra-stress tensor can be expressed as (see \\eqref{weak-est})\n$$\\TT := k \\left(\\sum_{i=1}^K \\CC_i(\\psi)\\right)-\\left(k(K+1) \\eta + \\de\\, \\eta^2\\right) \\II. $$\nThanks to the bounds on $\\eta$ stated in \\eqref{est-r-eta1} and \\eqref{est-r-eta3}, we have that\n\\be\\label{est-TT1}\n\\eta^2\\in L^\\infty(0,T;L^1(\\O)) \\cap L^1(0,T;L^3(\\O)).\n\\ee\n\nAccording to \\eqref{pt1-Mi-Ui}, we have $M=0$ on $\\d D$. We then deduce by using \\eqref{def-Ci} and \\eqref{pt2-M} that\n\\ba\\label{est-TT2}\n\\CC_i(\\psi) &= \\CC_i(M\\tpsi) = \\int_D M \\tpsi\\, U_i'\\bigg(\\frac{q_i^2}{2}\\bigg)\\, q_i q_i^{\\rm T} \\, \\dq = - \\int_{D} \\tpsi\\, (\\nabla_{q_i} M)q_i^{\\rm T}\\, \\dq\\\\\n&=\\int_{D} M\\, (\\nabla_{q_i}\\tpsi) q_i^{\\rm T}\\, \\dq + \\left(\\int_{D} M \\tpsi \\, \\dq\\right)\\II = \\int_{D} M (\\nabla_{q_i}\\tpsi) q_i^{\\rm T}\\, \\dq + \\eta \\II.\n\\ea\nThus, by \\eqref{nabla-psi}, H\\\"older's inequality implies that\n\\ba\\label{est-TT4}\n\\left|\\int_{D} M (\\nabla_{q_i}\\tpsi) q_i^{\\rm T}\\, \\dq \\right| \\leq c \\left(\\int_{D} M \\left(\\nabla_{q_i}\\sqrt{\\tpsi}\\right)^2 \\, \\dq \\right)^{\\frac 12}\\left(\\int_{D} M \\tpsi \\, \\dq \\right)^{\\frac 12}= c \\left(\\int_{D} M \\left(\\nabla_{q_i}\\sqrt{\\tpsi}\\right)^2 \\, \\dq \\right)^{\\frac 12} \\eta ^{\\frac 12}.\\nn\n\\ea\nBy \\eqref{est-psi1} we have\n$$\n\\left(\\int_{D} M \\left(\\nabla_{q_i}\\sqrt{\\tpsi}\\right)^2 \\dq \\right)^{\\frac 12} \\in L^2(0,T;L^2(\\O)).\n$$\nThus, by \\eqref{est-r-eta3} and \\eqref{est-TT1}, we obtain\n\\be\\label{est-TTf}\n\\TT \\in L^2(0,T;L^{\\frac 43}(\\O;\\R^{3\\times 3}))\\cap L^{\\frac 43}(0,T;L^\\frac{12}{7}(\\O;\\R^{3\\times 3})).\n\\ee\n\n\\subsection{Higher integrability of the fluid density and the pressure}\\label{sec:est-higher-vr}\n\nFrom the energy inequality \\eqref{energy} we deduce that\n$\\vr\\in L^\\infty(0,T;L^\\g(\\O))$, which implies that $p(\\vr)\\in L^\\infty(0,T;L^1(\\O))$; unfortunately, $p(\\vr)$ is only in $L^1$ with respect to the spatial variable, $x$. In order to improve the\nintegrability of the pressure with respect to the spatial variable, one may use the so-called Bogovski{\\u{\\i}} operator (see \\cite{bog}), exactly as in \\cite{FNP, F-book, N-book}.\nWe recall the following lemma whose proof can be found in Chapter III of Galdi's book \\cite{Galdi-book}.\n\\begin{lemma}\\label{lem-bog}\n Let $p \\in (1,\\infty)$ and suppose that $\\O\\subset \\R^d$ is a bounded Lipschitz domain. Let $L^{p}_0(\\Omega)$ be the space of $L^p(\\Omega)$ functions with zero mean-value.\nThen, there exists a linear operator $\\mathcal{B}_\\O$ from $L_0^p(\\O)$ to $W_0^{1,p}(\\O;\\R^d)$ such that\n$$\n\\Div_x\\, \\mathcal{B}_\\O (f) =f \\quad \\mbox{in} \\ \\O; \\quad \\|\\mathcal{B}_\\O (f)\\|_{W_0^{1,p}(\\O;\\R^d)} \\leq c (d,p,\\O)\\, \\|f\\|_{L^p(\\O)}\\quad\\, \\forall\\, f\\in L_0^p(\\O),\n$$\nwhere the constant $c$ depends only on $p$, $d$ and the Lipschitz character of $\\O$. If, in addition, $f=\\Div_x \\vg$ for some $\\vg \\in L^{q}, \\ 10$ sufficiently small, that\n\\be\\label{est-bog-p1}\n\\int_0^t\\int_\\O \\phi(t)\\, \\vr^\\g\\, b_{n}(\\vr)\\, \\dx \\, \\dt \\leq c\\qquad \\forall\\,n \\geq 0.\n\\ee\nLetting $n\\to \\infty$ in \\eqref{est-bog-p1} and approximating $1$ by $C_c^\\infty ((0,T))$ functions finally gives\n\\be\\label{est-bog-pf1}\n \\vr \\in L^{\\g+\\th}((0,T)\\times \\O),\\quad p(\\vr) \\in L^{1+\\frac{\\th}{\\g}}((0,T)\\times \\O) \\quad \\mbox{for some $\\th>0$}.\n\\ee\n\n\\medskip\n\n\n\\begin{remark}\\label{rem-bog-th}\n\n\n\nIn order to obtain \\eqref{est-bog-pf1} for some positive $\\th$, the conditions imposed in \\eqref{mu-eta2} and \\eqref{mu-eta3} can be relaxed. By careful analysis the following constraints are found to be sufficient:\n\\be\\label{mu-eta20}\n\\g>\\frac 32,\\ 0\\leq \\o < \\frac{10}{3} \\quad \\mbox{or} \\quad -2<\\o\\leq 0,\\ \\g>\\frac{6}{4+3\\o}.\n\\ee\nThe more restrictive conditions featuring in \\eqref{mu-eta2} and \\eqref{mu-eta3} are needed later on, in Section \\ref{vis-flux}, in order to prove the so-called {\\rm effective viscous flux} equality (see Remark \\ref{rem-mu-eta}).\n\\end{remark}\n\n\\subsection{Bounds on the time derivative and continuity}\\label{sec:est-time}\n\nThis section is devoted to establishing bounds on the time derivatives $(\\d_t \\vr, \\ \\d_t\\eta, \\ \\d_t (\\vr \\vu), \\ \\d_t \\psi)$.\n\nAs $\\d_t \\vr = - \\Div_x(\\vr\\vu)$, the bound in \\eqref{est-ru2} implies that\n$\\d_t \\vr \\in L^\\infty(0,T; W^{-1,\\frac{2\\g}{\\g+1}}(\\O)).$ Then, by \\eqref{est-r-eta1} and Lemma \\ref{lem-Cw-Cw*}, we have that\n$$\n\\vr \\in C_w([0,T]; L^\\g(\\O)).\n$$\n\n\n\\medskip\n\nRecall that $\\d_t \\eta = - \\Div_x(\\eta\\vu) + \\e \\Delta_x \\eta$. For $\\o\\geq 0$ satisfying \\eqref{mu-eta2}, we have $\\vu \\in L^2(0,T; W_0^{1,2}(\\O;\\R^3))$, which is embedded in $L^2(0,T; L^6(\\O;\\R^3))$. Thanks to \\eqref{est-r-eta1}, we then obtain $\\d_t \\eta \\in L^2(0,T; W^{-1,\\frac{3}{2}}(\\O)).$\nIf on the other hand $\\o\\leq 0$ satisfies \\eqref{mu-eta3}, then we have $\\vu \\in L^2(0,T;W_0^{1,\\frac{4}{2+|\\o|}}(\\O;\\R^3)) \\hookrightarrow L^2(0,T; L^\\frac{12}{3|\\o|+2}(\\O;\\R^3))$. Consequently, $\\d_t \\eta \\in L^2(0,T; W^{-1,\\frac{12}{3|\\o|+8}}(\\O))$.\nThus, in both cases, by \\eqref{est-r-eta1} and Lemma \\ref{lem-Cw-Cw*}, we have\n$$\n\\eta \\in C_w([0,T]; L^2(\\O)).\n$$\n\n\\medskip\n\nNext, we recall that $\\d_t(\\vr\\vu)=- \\Div_x (\\vr \\vu \\otimes \\vu) - \\nabla_x p(\\vr) +\n\\Div_x \\SSS +\\Div_x \\TT +\\vr\\, \\ff$. We first consider the case with $\\o\\geq 0$ satisfying \\eqref{mu-eta2}. By the estimates in Section \\ref{sec:est-u} we have $\\vu \\in L^2(0,T; L^6(\\O;\\R^3))$. Since $\\vr\\in L^\\infty(0,T;L^\\g(\\O))$ with $\\g>\\frac{3}{2}$, together with \\eqref{est-r-eta1}, we obtain\n\\be\\label{est-ruu}\n\\vr \\vu \\otimes \\vu \\in L^1(0,T; L^{\\frac{3\\g}{\\g+3}}(\\O;\\R^{3\\times 3}))\\cap L^\\infty(0,T; L^{1}(\\O;\\R^{3\\times 3})) \\hookrightarrow L^r(0,T; L^{r}(\\O;\\R^{3\\times 3})) \\quad \\mbox{for some $r>1$}.\n\\ee\nBy \\eqref{est-SS1}, \\eqref{est-bog-pf1}, \\eqref{est-TTf}, together with \\eqref{est-ruu}, we deduce that\n\\be\\label{est-dt-ru}\n\\d_t (\\vr\\vu) \\in L^r(0,T; W^{-1,r}(\\O;\\R^{3})) \\quad \\mbox{for some $r>1$}.\n\\ee\nFor the case with $\\o\\leq 0$ satisfying \\eqref{mu-eta3}, a similar argument gives the same result as in \\eqref{est-dt-ru}. Then, by \\eqref{est-ru2}, \\eqref{est-dt-ru} and Lemma \\ref{lem-Cw-Cw*}, we have in both cases that\n$$\\vr\\vu \\in C_w([0,T]; L^\\frac{2\\g}{\\g+1}(\\O;\\R^3)).$$\n\n\nBy the Fokker--Planck equation \\eqref{eq-psi} for $\\psi$, the estimates in \\eqref{est-psi9} and \\eqref{est-psi10}, direct calculations yield that\n\\be\\label{est-dt-psi}\n\\d_t \\psi \\in\nL^2(0,T; W^{s,2}(\\O\\times D)') \\quad \\mbox{with $s>1 + \\frac{3}{2}(K+1)$},\n\\ee\nwhere we have used that $W^{s,2}(\\Omega \\times D) \\hookrightarrow W^{1,\\infty}(\\Omega \\times D)$ for $s>1+\\frac{3}{2}(K+1)$.\nMoreover, by the estimates \\eqref{est-psi5} and \\eqref{est-dt-psi}, using Lemma \\ref{lem-Cw-Cw*} and the same argument as in Section 4.5 in \\cite{Barrett-Suli}, we deduce that $$\\psi \\in C_w([0,T];L^1(\\O\\times D)).$$\n\n\n\n\n\\subsection{Summary}\n\nWe summarize the results obtained in this section. Let $(\\vr,\\vu,\\psi,\\eta)$ be a dissipative weak solution in the sense of Definition \\ref{def-weaksl} associated\nwith initial data $(\\vr_0,\\vu_0,\\psi_0,\\eta_0)$ satisfying \\eqref{ini-data}. Then,\n\\ba\\label{est-sum1}\n& \\vr\\in C_w([0,T);L^\\gamma(\\O))\\cap L^{\\g+\\th}((0,T)\\times \\O);\\quad \\eta \\in C_w([0,T];L^2(\\O)) \\cap L^2(0,T;W^{1,2}(\\O));\\\\\n& \\vr |\\vu|^2 \\in L^{\\infty}(0,T;L^{1}(\\O)),\\quad \\vr \\vu \\in C_w([0,T];L^{\\frac{2\\g}{\\g+1}}(\\O;\\R^3));\\quad \\TT \\in L^2(0,T;L^{\\frac 43}(\\O;\\R^{3\\times 3}))\\cap L^{\\frac 43}(0,T;L^\\frac{12}{7}(\\O;\\R^{3\\times 3}));\\\\\n&\\vu \\in L^2(0,T;W^{1,2}_0(\\O;\\R^3)) \\ \\mbox{if $\\o\\geq 0$ satisfies \\eqref{mu-eta2}};\\\\\n& \\vu \\in L^2(0,T;W_0^{1,\\frac{4}{2+|\\o|}}(\\O;\\R^{3})) \\cap L^{\\frac{4}{2+|\\o|}}(0,T;W_0^{1,\\frac{12}{6+|\\o|}}(\\O;\\R^{3}))\\cap L^{\\frac{20}{10+3|\\o|}}(0,T;W_0^{1,\\frac{20}{10+3|\\o|}}(\\O;\\R^{3})) \\\\\n&\\qquad \\mbox{if $\\o\\leq 0$ satisfies \\eqref{mu-eta3}};\\\\\n&\\mathcal{F}(\\tpsi), \\tpsi\\log\\tpsi, \\ \\tpsi \\in L^\\infty(0,T;L^1_M(\\O\\times D)),\\quad \\sqrt{\\tpsi} \\in L^2(0,T;H_M^1(\\O\\times D)),\\\\\n&\\qquad \\nabla_{q,x} \\tpsi \\in L^2(0,T;L^{\\frac 43}(\\O;L^1_M(D;\\R^{3(K+1)})))\\cap L^{\\frac 43}(0,T;L^\\frac{12}{7}(\\O;L^1_M(D;\\R^{3(K+1)})));\\\\\n&\\mathcal{F}(\\psi),\\ \\psi\\log\\psi \\ \\in L^\\infty(0,T;L^1(\\O\\times D)),\\quad \\psi \\in C_w([0,T];L^1(\\O\\times D)),\\\\\n &\\qquad \\nabla_x \\psi \\in L^2(0,T;L^{\\frac 43}(\\O;L^1(D;\\R^{3})))\\cap L^{\\frac 43}(0,T;L^\\frac{12}{7}(\\O;L^1(D;\\R^{3})));\n\\ea\nand\n\n\\begin{equation}\\label{est-sum2}\n\\begin{aligned}\n&\\d_t \\vr \\in L^\\infty(0,T; W^{-1,\\frac{2\\g}{\\g+1}}(\\O)), \\quad \\d_t (\\vr\\vu) \\in L^r(0,T; W^{-1,r}(\\O;\\R^{3})) &&\\quad \\mbox{for some $r>1$};\\\\\n& \\d_t \\eta \\in L^2(0,T; W^{-1,\\frac{3}{2}}(\\O)) &&\\quad \\mbox{if $\\o\\geq 0$ satisfies \\eqref{mu-eta2}};\\\\\n&\\d_t \\eta \\in L^2(0,T; W^{-1,\\frac{12}{3|\\o|+8}}(\\O)) &&\\quad \\mbox{if $\\o\\leq 0$ satisfies \\eqref{mu-eta3}};\\\\\n&\\d_t \\psi \\in L^2(0,T; W^{s,2}(\\O\\times D)') &&\\quad \\mbox{with $s>1 + \\frac{3}{2}(K+1)$}.\n\\end{aligned}\n\\end{equation}\nIt is important to note that all of the above inclusions are consequences of bounds that depend only on the initial energy, the final time $T$, and the structural constants in the hypotheses imposed on the constitutive relations.\n\n\n\n\\section{Passing to the limit}\\label{sec:lim}\n\nLet $(\\vr_n,\\vu_n,\\psi_n,\\eta_n)_{n\\in \\N}$ be a sequence of the dissipative (finite-energy) weak solutions satisfying the assumptions in Theorem \\ref{theorem}. Then, the energy inequality \\eqref{energy} and \\eqref{energy-f} give, respectively,\n\\ba\\label{energy-n}\n&\\int_{\\Omega} \\bigg[\n\\frac{1}{2} \\vr_n |\\vu_n|^2 + P(\\vr_n) +\\de\\, \\eta_n^2 + k \\int_D M \\mathcal{F} (\\psi_n)\\, \\dq\\bigg] (t, \\cdot) \\, \\dx + 2\\, \\e\\, \\de \\int_0^t \\intO{ |\\nabla_x \\eta_n|^2} \\, \\dt'\\\\\n&\\quad+ \\int_0^t \\intO{ \\mu^S(\\eta_n) \\left|\\frac{\\nabla \\vu_n + \\nabla^{\\rm T} \\vu_n}{2} - \\frac{1}{3} (\\Div_x \\vu_n) \\II \\right|^2 +\\mu^B(\\eta_n) |\\Div_x \\vu_n|^2 } \\, \\dt' \\\\\n&\\quad+ \\e\\, k \\int_0^t \\intO{ \\int_D M\\left|\\nabla_x \\sqrt{\\tpsi_n} \\right|^2 \\dq} \\, \\dt+\\frac{k\\,A_0}{4\\lambda} \\int_0^t \\intO{ \\int_D M \\left|\\nabla_{q} \\sqrt{\\tpsi_n}\\right|^2 \\dq} \\, \\dt'\\\\\n&\\leq \\int_{\\Omega} \\bigg[\n\\frac{1}{2} \\vr_{0,n} |\\vu_{0,n} |^2 + P(\\vr_{0,n}) + \\de\\, \\eta_{0,n}^2 + k \\int_D M \\mathcal{F} \\bigg(\\frac{\\psi_{0,n}}{M}\\bigg)\\,\\dq\\bigg]\\dx + \\int_0^t\\int_\\O \\vr_n \\,\\ff \\cdot \\vu_n \\, \\dx\\, \\dt'\n\\ea\n and, therefore,\n \\ba\\label{energy-f-n}\n&\\int_{\\Omega} \\bigg[ \\frac{1}{2} \\vr_n |\\vu_n|^2 + P(\\vr_n) +\\de\\, \\eta_n^2 + k \\int_D M \\mathcal{F} (\\widetilde \\psi_n)\\,\\dq\\bigg] (t, \\cdot) \\, \\dx + 2\\, \\e\\, \\de \\int_0^t \\intO{ |\\nabla_x \\eta_n|^2} \\, \\dt' \\\\\n&\\quad+ \\int_0^t \\intO{ \\mu^S(\\eta_n) \\left|\\frac{\\nabla \\vu_n + \\nabla^{\\rm T} \\vu_n}{2} - \\frac{1}{3} (\\Div_x \\vu_n) \\II \\right|^2 +\\mu^B(\\eta_n) |\\Div_x \\vu_n|^2 } \\ \\dt' \\\\\n&\\quad+\\e\\, k \\int_0^t \\intO{ \\int_D M\\left|\\nabla_x \\sqrt{\\widetilde \\psi_n} \\right|^2\\dq} \\ \\dt'+\\frac{k\\,A_0}{4\\lambda} \\int_0^t \\intO{ \\int_D M \\left|\\nabla_{q} \\sqrt{\\widetilde \\psi_n}\\right|^2\\dq} \\, \\dt' \\\\\n&\\leq c(1+T)\\, {\\rm e}^t.\n\\ea\nThus, from the results in Section \\ref{sec:estimate}, this solution sequence $(\\vr_n,\\vu_n,\\psi_n,\\eta_n)_{n\\in \\N}$ satisfies the uniform bounds \\eqref{est-sum1} and \\eqref{est-sum2}.\n\nThe present section is devoted to studying the limit of this solution sequence. We remark that, throughout this section, the limits are taken up to subtractions of subsequences without identification.\n\n\n\\subsection{Convergence of the fluid density}\\label{sec:vr-lim}\n\nWe have that $\\{\\vr_n\\}_{n\\in \\N}$ is a sequence in $C_w([0,T];L^\\gamma(\\O))$ satisfying\n\\be\\label{vrn-est}\n\\sup_{n\\in \\N} \\left(\\|\\vr_n\\|_{L^\\infty(0,T;L^\\gamma(\\O))} + \\|\\d_t \\vr_n \\|_{L^\\infty(0,T; W^{-1,\\frac{2\\g}{\\g+1}}(\\O))} \\right) \\leq c.\n\\ee\nBy Sobolev embedding one has\n$$\nL^\\g(\\O) \\hookrightarrow W^{-1,\\frac{3\\g}{3-\\g}}(\\Omega) \\quad \\mbox{ if $\\g<3$};\\quad L^\\g(\\O) \\hookrightarrow W^{-1,s}(\\Omega) \\quad \\mbox{ for any $s \\in (1,\\infty)$, \\ if $\\g\\geq3$}.\n$$\nThus, by applying Lemma \\ref{lem-Cw}, we deduce that\n\\ba\\label{vrn-lim}\n& \\vr_n \\to \\vr \\quad \\mbox{in}\\ C_w([0,T];L^\\gamma(\\O));\\\\\n& \\vr_n \\to \\vr \\quad \\mbox{strongly in} \\ C([0,T];W^{-1,r}(\\O)) \\quad \\mbox{for any $r\\geq \\frac{3}{2}$ such that $\\frac{3r}{3+r}<\\g$}.\n\\ea\n\n\\subsection{Convergence of the fluid velocity field}\\label{sec:lim-u}\n\nBy Sobolev embedding, for the case $\\o\\geq 0$, we have\n\\be\\label{vun-limit1}\n\\vu_n \\to \\vu \\quad \\mbox{weakly in}\\ L^2(0,T;W^{1,2}_0(\\O;\\R^3))\\ \\mbox{and in}\\ L^2(0,T;L^{6}(\\O;\\R^3)),\n\\ee\nwhile for the case $\\o\\leq 0$, we have\n\\ba\\label{vun-limit2}\n&\\vu_n \\to \\vu \\quad \\mbox{weakly in}\\ L^2(0,T;W_0^{1,\\frac{4}{2+|\\o|}}(\\O;\\R^3))\\ \\mbox{and weakly in}\\ L^2(0,T;L^{\\frac{12}{2+3|\\o|}}(\\O;\\R^3)),\\\\\n&\\vu_n \\to \\vu \\quad \\mbox{weakly in}\\ L^{\\frac{4}{2+|\\o|}}(0,T;W_0^{1,\\frac{12}{6+|\\o|}}(\\O;\\R^3))\\ \\mbox{and weakly in}\\ L^{\\frac{4}{2+|\\o|}}(0,T;L^{\\frac{12}{2+|\\o|}}(\\O;\\R^3)),\\\\\n&\\vu_n \\to \\vu \\quad \\mbox{weakly in}\\ L^{\\frac{20}{10+3|\\o|}}(0,T;W_0^{1,\\frac{20}{10+3|\\o|}}(\\O;\\R^3)) \\ \\mbox{and weakly in}\\ L^{\\frac{20}{10+3|\\o|}}(0,T;L^{\\frac{60}{10+9|\\o|}}(\\O;\\R^3)).\n\\ea\n\n\\subsection{Convergence of the nonlinear terms $\\vr_n\\vu_n$ and $\\vr_n \\vu_n\\otimes \\vu_n$}\\label{sec:lim-convective}\n\nWe first consider the case with $\\o\\geq 0$ satisfying \\eqref{mu-eta2}. Since $\\g>\\frac 32>\\frac 65$, we have, because of \\eqref{vrn-lim}, that\n\\ba\\label{vrn-lim1}\n \\vr_n \\to \\vr \\quad \\mbox{strongly in} \\ C([0,T];W^{-1,2}(\\O)).\n\\ea\nTogether with \\eqref{vun-limit1}, we have\n\\be\\label{vrvun-lim0}\n\\vr_n\\vu_n \\to \\vr\\vu \\quad \\mbox{in $\\mathcal{D}'((0,T)\\times \\O;\\R^{3})$.}\n\\ee\nBy observing the uniform estimate\n$$\n\\sup_{n \\in \\mathbb{N}}\\left(\\|\\vr_n\\vu_n\\|_{L^\\infty(0,T;L^{\\frac{2\\g}{1+\\g}}(\\O;\\R^{3}))} + \\|\\vr_n\\vu_n\\|_{L^2(0,T;L^{\\frac{6\\g}{6+\\g}}(\\O;\\R^{3}))}\\right)\\leq c\n$$\nwe have\n\\be\\label{vrvun-lim1}\n\\vr_n\\vu_n \\to \\vr\\vu \\quad \\mbox{weakly* in} \\ L^\\infty(0,T;L^{\\frac{2\\g}{1+\\g}}(\\O;\\R^{3})) \\ \\mbox{and weakly in} \\ L^2(0,T;L^{\\frac{6\\g}{6+\\g}}(\\O;\\R^{3})).\n\\ee\n\nMoreover, by \\eqref{est-sum1} and \\eqref{est-sum2}, we have $\\vr_n\\vu_n \\in C_w([0,T];L^{\\frac{2\\g}{\\g+1}}(\\O;\\R^d))$ and $\\{\\d_t (\\vr_n\\vu_n)\\}_{n\\in \\N}$ is uniformly bounded in $L^r(0,T;W^{-1,r}(\\O;\\R^3))$ for some $r>1$. By Lemma \\ref{lem-Cw} and the interpolation argument in Section \\ref{sec:vr-lim} we have\n\\ba\\label{vrvun-lim2}\n&\\vr_n\\vu_n \\to \\vr\\vu \\quad \\mbox{in}\\ C_w([0,T];L^{\\frac{2\\g}{1+\\g}}(\\O;\\R^3));\\\\\n&\\vr_n\\vu_n \\to \\vr\\vu \\quad \\mbox{strongly in} \\ C([0,T];W^{-1,r}(\\O;\\R^3)) \\quad \\mbox{for any $r\\geq \\frac{3}{2}$ such that $\\frac{3r}{3+r}<\\frac{2\\g}{1+\\g}$}.\n\\ea\n\nThe fact that $\\g>\\frac 32$ implies $\\frac{2\\g}{1+\\g} > \\frac{6}{5}$. This gives $\\vr_n\\vu_n \\to \\vr\\vu$ strongly in $C([0,T];W^{-1,2}(\\O;\\R^3))$. Together with \\eqref{vun-limit1}, we have for the sequence of convective terms:\n\\be\\label{vrvuvun-lim0}\n\\vr_n\\vu_n\\otimes \\vu_n \\to \\vr\\vu\\otimes \\vu \\quad \\mbox{ in $\\mathcal{D}'((0,T)\\times \\O;\\R^{3\\times 3})$.}\n\\ee\n\n\nThe case $\\o\\leq 0$ is dealt with similarly to the case $\\o\\geq 0$, and we obtain the results stated in \\eqref{vrvun-lim0} and \\eqref{vrvuvun-lim0}, so we omit the details. We merely remark that when $\\o\\leq 0$, then we have weaker integrability for $\\vu$; this is, however, compensated by supposing stronger integrability for $\\vr$ through hypothesis \\eqref{mu-eta3}.\n\n We also remark that the hypotheses stated in \\eqref{mu-eta2} and \\eqref{mu-eta3} can be relaxed in this part of the analysis: the constraints stated in \\eqref{mu-eta20} in Remark \\ref{rem-bog-th} are sufficient.\n\n\\subsection{Convergence of the polymer number density}\\label{sec:lim-eta}\n\nWe begin by noting that the sequence $\\{\\eta_n\\}_{n\\in \\N}$ is contained in $C_w([0,T];L^2(\\O))$ and satisfies the bound\n\\be\\label{etan-est1}\n\\sup_{n \\in \\mathbb{N}} \\left(\\|\\eta_n\\|_{L^\\infty(0,T;L^2(\\O)) \\cap L^2(0,T;W^{1,2}(\\O))} + \\|\\d_t \\eta_n\\|_{L^2(0,T;W^{-1,\\frac{12}{3|\\o|+8}}(\\O))}\\right) \\leq c.\n\\ee\nThus, by Lemma \\ref{lem-Cw},\n\\be\\label{etan-lim1}\n\\eta_n \\to \\eta \\quad \\mbox{strongly in}\\ C_w([0,T];L^2(\\O)) \\ \\mbox{and weakly in}\\ L^2(0,T;W^{1,2}(\\O)).\n\\ee\nThanks to the compact Sobolev embedding $W^{1,2}(\\O) \\hookrightarrow \\hookrightarrow L^q(\\O)$, $q<6$, the\nAubin--Lions--Simon compactness theorem (see \\cite{Lions-InC} or \\cite{Simon}) implies that\n\\be\\label{etan-lim2}\n\\eta_n \\to \\eta \\quad \\mbox{strongly in}\\ L^2(0,T;L^q(\\O))\\quad \\mbox{for any $q<6$}.\n\\ee\nBy interpolation we also have\n\\be\\label{etan-lim3}\n\\eta_n \\to \\eta \\quad \\mbox{strongly in}\\ L^q((0,T) \\times \\O)\\quad \\mbox{for any $q<\\frac{10}{3}$}.\n\\ee\n\nWe still need to show that $\\eta=\\int_D \\psi \\ \\dq$, where $\\psi$ is the limit of $\\psi_n$. This will be done later on, in Section \\ref{sec:limit-psi}.\n\n\n\\subsection{Convergence of the Newtonian stress tensor}\\label{sec:lim-SS}\nWhen $0\\leq \\o <\\frac{5}{3}$ as in \\eqref{mu-eta2}, by \\eqref{mu-eta1} and \\eqref{etan-lim3}, we have that\n\\be\\label{mu-lim1}\n(\\mu^S(\\eta_n),\\mu^B(\\eta_n)) \\to (\\mu^S(\\eta),\\mu^B(\\eta)) \\quad \\mbox{strongly in}\\ L^{\\frac{10}{3\\o}}((0,T)\\times \\O).\n\\ee\nTogether with \\eqref{vun-limit1}, we deduce that\n\\be\\label{Newtonian-lim1}\n\\SSS_n \\to \\SSS :=\\mu^S(\\eta) \\left( \\frac{\\nabla \\vu + \\nabla^{\\rm T} \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu) \\II \\right) + \\mu^B(\\eta) (\\Div_x \\vu) \\II\\quad {\\rm weakly\\ in} \\ L^{\\frac{10}{3\\o+5}}((0,T)\\times \\O;\\R^{3\\times 3}).\n\\ee\n\nOn the other hand, when $-\\frac{4}{3}<\\o\\leq 0$ as in \\eqref{mu-eta2}, then we have\n\\be\\label{mu-lim2}\n(\\mu^S(\\eta_n),\\mu^B(\\eta_n)) \\to (\\mu^S(\\eta),\\mu^B(\\eta)) \\quad \\mbox{strongly in}\\ L^q(0,T;L^q(\\O)) \\quad \\mbox{for any $q<\\infty$}.\n\\ee\nThus, by \\eqref{vun-limit2}, we have, for any $r<\\frac{20}{3\\o+10}$,\n\\ba\\label{Newtonian-lim10}\n\\SSS_n \\to \\SSS :=\\mu^S(\\eta) \\left( \\frac{\\nabla \\vu + \\nabla^{\\rm T} \\vu}{2} - \\frac{1}{3} (\\Div_x \\vu) \\II \\right) + \\mu^B(\\eta) (\\Div_x \\vu) \\II\\quad\n{\\rm weakly~in} \\ L^{r}((0,T)\\times \\O;\\R^{3\\times 3}).\n\\ea\n\n\nWe remark that, by rewriting\n\\ba\\label{Newtonian-est2}\n\\SSS_n = \\sqrt{\\mu^S(\\eta_n)} \\sqrt{\\mu^S(\\eta_n)} \\left( \\frac{\\nabla \\vu_n + \\nabla^{\\rm T} \\vu_n}{2} - \\frac{1}{3} (\\Div_x \\vu_n) \\II \\right) + \\sqrt{\\mu^B(\\eta_n)}\\sqrt{\\mu^B(\\eta_n)} (\\Div_x \\vu_n) \\II,\n\\ea\nwe can relax the constraints in \\eqref{mu-eta2} and \\eqref{mu-eta3} as in \\eqref{mu-eta20} in Remark \\ref{rem-bog-th}.\n\n\\subsection{Convergence of the probability density function}\\label{sec:limit-psi}\n\nFirst of all, by the energy inequality \\eqref{energy-f-n} we have\n\\be\\label{est-psin1}\n\\sup_{n\\in\\N}\\left( \\left\\|\\mathcal{F}(\\tpsi_n)\\right\\|_{L^\\infty(0,T;L^1_M(\\O\\times D))} + \\left\\|\\sqrt{\\tpsi_n}\\,\\right\\|_{L^2(0,T;H_M^1(\\O\\times D))} \\right)\\leq c.\n\\ee\n\n\nAs in Section 4 in \\cite{Barrett-Suli} and Section 5 in \\cite{Barrett-Suli4}, we use Dubinski{\\u\\i}'s compactness theorem (cf. \\cite{Dubin}; see also \\cite{Barrett-Suli3}, Theorem 3.1 or \\cite{Barrett-Suli}) by setting\n\\ba\\label{dubinskii-spaces}\n&X:= L^1_M(\\O\\times D),\\quad X_0:=\\{ \\vp\\in X : \\vp\\geq 0,\\ \\sqrt{\\vp} \\in H^1_M(\\O\\times D) \\},\\\\\n&X_1: = M^{-1} W^{s,2}(\\O\\times D)':=\\{ M^{-1} \\vp : \\vp \\in W^{s,2}(\\O\\times D)'\\} \\quad \\mbox{with $s>1+\\frac{3}{2}(K+1)$},\n\\ea\nwhere $X_0$ is a seminomed space (in the sense of Dubinski{\\u\\i}) with seminorm defined by\n$$\n[\\vp]_{X_0}: = \\|\\vp\\|_{X} + \\int_{\\O\\times D} M\\bigg( \\left|\\nabla_x \\sqrt{\\vp}\\,\\right|^2 + \\left|\\nabla_q \\sqrt{\\vp}\\,\\right|^2 \\bigg) \\, \\dq \\, \\dx.\n$$\n\nBy \\eqref{est-sum1} and \\eqref{est-sum2} we have that\n\\be\\label{psin-est1}\n\\mbox{$\\{\\widetilde \\psi_n\\}_{n\\in \\N}$ is uniformly bounded in $L^1(0,T;X_0)$\\quad and \\quad $\\{\\d_t \\widetilde \\psi_n\\}_{n\\in \\N}$ is uniformly bounded in $L^2(0,T;X_1)$}.\n\\ee\nThe continuity of the embedding $X\\hookrightarrow X_1$ and the compactness of the embedding $X_0 \\hookrightarrow X$ are shown in Section 5 in \\cite{Barrett-Suli4}. Then, by virtue of Dubinski{\\u\\i}'s compact embedding theorem, we have that\n\\be\\label{psin-lim1}\n\\tpsi_n \\to \\tpsi \\quad \\mbox{strongly in} \\ L^1(0,T;L_M^1(\\Omega\\times D)),\n\\ee\nwhich is equivalent to\n\\be\\label{psin-lim2}\n\\psi_n \\to \\psi \\quad \\mbox{strongly in} \\ L^1(0,T;L^1(\\Omega\\times D)).\n\\ee\nThis implies that\n\\be\\label{psin-lim3}\n\\eta_n = \\int_D \\psi_n \\,\\dq \\to \\int_D \\psi \\, \\dq \\quad \\mbox{strongly in} \\ L^1((0,T)\\times \\Omega);\n\\ee\nin addition, by the uniqueness of the limit, the function $\\eta$ obtained in Section \\ref{sec:lim-eta} satisfies $\\eta = \\int_D \\psi \\, \\dq$.\n\n\n\nSince $\\mathcal{F}(\\tpsi)$ is nonnegative on $(0,T)\\times \\Omega \\times D$, by applying Fatou's lemma we have that\n\\be\\label{psin-est3}\n \\left\\|\\mathcal{F}(\\tpsi)\\right\\|_{L^\\infty(0,T;L^1_M(\\O\\times D))} \\leq \\liminf_{n\\to \\infty} \\left\\|\\mathcal{F}(\\tpsi_n)\\right\\|_{L^\\infty(0,T;L^1_M(\\O\\times D))} \\leq c.\n\\ee\nBy the same technique as in Section \\ref{sec:est-psi} we deduce from \\eqref{psin-est3} that\n\\be\\label{psin-est4}\n\\tpsi\\log\\tpsi, \\ \\tpsi \\in L^\\infty(0,T;L^1_M(\\O\\times D)),\\quad \\psi\\log\\psi, \\ \\psi \\in L^\\infty(0,T;L^1(\\O\\times D)).\n\\ee\n\nBy \\eqref{psin-est1} we have\n\\be\\label{psin-est5}\n\\d_t \\psi \\in L^2(0,T;W^{s,2}(\\O\\times D)') \\quad \\mbox{for any $s>1 + \\frac{3}{2}(K+1)$}.\n\\ee\n\nFrom the estimates in \\eqref{psin-est4} and \\eqref{psin-est5}, by using the second part of Lemma \\ref{lem-Cw-Cw*} and the same argument as in Section 4.5 in \\cite{Barrett-Suli}, we deduce that\n$\n \\psi \\in C_w([0,T];L^1(\\O\\times D)).\n$\n\n\nAgain, we write $\\nabla_{q_i}\\tpsi_n = 2 \\sqrt{\\tpsi_n}\\nabla_{q_i}\\sqrt{\\tpsi_n}$. By \\eqref{psin-lim1} we have\n$$\n\\sqrt{\\tpsi_n} \\to \\sqrt{\\tpsi} \\quad \\mbox{strongly in} \\ L^2(0,T;L_M^2(\\Omega\\times D)).\n$$\nFurther, by \\eqref{est-psin1}, we have that\n$$\n\\nabla_q \\sqrt{\\tpsi_n} \\to \\nabla_q \\sqrt{\\tpsi} \\quad \\mbox{weakly in} \\ L^2(0,T;L_M^2(\\Omega\\times D;\\R^{3K})).\n$$\nThus,\n\\be\\label{psin-limf1}\n\\nabla_{q}\\tpsi_n \\to \\nabla_{q}\\tpsi \\quad \\mbox{weakly in} \\ L^1(0,T;L_M^1(\\Omega\\times D;\\R^{3K})).\n\\ee\n\n\nSimilarly, we also have that\n\\be\\label{psin-limf2}\n\\nabla_{x}\\tpsi_n \\to \\nabla_{x}\\tpsi \\quad \\mbox{weakly in} \\ L^1(0,T;L_M^1(\\Omega\\times D;\\R^{3})),\\quad\n\\nabla_{x}\\psi_n \\to \\nabla_{x}\\psi \\quad \\mbox{weakly in} \\ L^1(0,T;L^1(\\Omega\\times D;\\R^{3})).\n\\ee\n\n\n\\subsection{Convergence of the extra-stress tensor}\\label{sec:lim-TT}\n\nWe recall the formula for the extra-stress tensor. By \\eqref{weak-est} and \\eqref{est-TT2} we have that\n\\be\\label{TTn-est1}\n\\TT_n := k \\left(\\sum_{i=1}^K \\int_{D} M (\\nabla_{q_i}\\tpsi_n) q_i^{\\rm T}\\, \\dq \\right) - \\left(k \\,\\eta_n + \\de\\, \\eta_n^2\\right) \\II,\n \\ee\nwhich, because of \\eqref{est-TTf}, satisfies\n\\be\\label{TTn-est2}\n\\sup_{n\\in\\N}\\left(\\|\\TT_n\\|_{L^2(0,T;L^{\\frac 43}(\\O;\\R^{3\\times 3}))} + \\|\\TT_n\\|_{L^{\\frac 43}(0,T;L^\\frac{12}{7}(\\O;\\R^{3\\times 3}))}\\right) -\\frac{4}{3}$, appearing in \\eqref{mu-eta3}, is really needed. Indeed, for any test function $\\vvarphi$, just as in \\eqref{nonlin-FP-lim3}, we have that\n\\be\\label{nonlin-FP-lim5}\n\\int_D (\\psi_n-\\psi) \\, \\vvarphi \\,q_i \\ \\dq \\to \\mathbf{0} \\quad \\mbox{in}\\ L^q(0,T;L^q(\\O;\\R^3)) \\quad \\mbox{for any $q<\\frac{10}{3}$}.\n\\ee\nAccording to the bounds established in Section \\ref{sec:lim-u} for the case $\\o\\leq 0$, the sequence $\\left\\{\\nabla_x \\vu_n\\right\\}_{n\\in \\N}$ is uniformly bounded in $L^{\\frac{20}{10+3|\\o|}}((0,T)\\times\\O;\\R^{3\\times 3})$. To deduce the desired convergence result, we therefore need that\n\\be\\label{mu-eta30}\n\\frac{10+3|\\o|}{20} + \\frac{3}{10}<1 \\Longleftrightarrow |\\o| < \\frac{4}{3}.\n\\ee\n\n\n\\subsection{Convergence of the fluid pressure}\nAs was shown in Section \\ref{sec:est-higher-vr} concerning higher integrability of the fluid density and pressure, we have that\n\\be\\label{pn-lim1}\np(\\vr_n) \\to \\overline{p(\\vr)} \\quad \\mbox{weakly in} \\ L^{1+\\frac{\\th}{\\g}}((0,T)\\times \\O).\n\\ee\n\nBy applying this, together with the convergence results from the previous sections, in particular the convergence results for the nonlinear terms stated in Section \\ref{sec:lim-convective}, Section \\ref{sec:lim-SS} and Section \\ref{sec:lim-TT}, we deduce that\n\\be\\label{i1-0}\n\\d_t \\vr + \\Div_x (\\vr \\vu) = 0 \\quad \\mbox{in}\\ \\mathcal{D}'((0,T)\\times \\O),\n\\ee\n\\be\\label{i2-0}\n \\d_t (\\vr\\vu)+ \\Div_x (\\vr \\vu \\otimes \\vu) +\\nabla_x \\overline{p(\\vr)} -\n\\Div_x \\SSS =\\Div_x \\TT +\\vr\\, \\ff\\quad \\mbox{in}\\ \\mathcal{D}'((0,T)\\times \\O;\\R^3),\n\\ee\nwhere the Newtonian stress tensor $\\SSS$ and the extra-stress tensor $\\TT$ are defined by \\eqref{Newtonian-lim1} and \\eqref{TTn-lim2}, respectively.\n\n\nConcerning the convergence of the nonlinear terms in the Navier--Stokes--Fokker--Planck system with variable viscosity coefficients considered here, it remains to show that $p(\\vr) = \\overline{p(\\vr)},$\nwhich, because of the strict convexity of $p(\\cdot)$, is equivalent to the strong convergence of $\\vr_n$:\n\\be\\label{pn-limf2}\n\\vr_n \\to \\vr \\quad \\mbox{a.e. in} \\ (0,T)\\times \\O.\n\\ee\n\nThis is also one of the main difficulties in the study of global existence of weak solutions to the compressible Navier--Stokes equations (see \\cite{Lions-C, FNP,F-book,N-book}), where the so called \\emph{effective viscous flux} introduced by P. L. Lions \\cite{Lions-C} plays a crucial role. It turns out that the effective viscous flux as a whole is more regular than its components.\n We will prove the following lemma, which is in the spirit of Proposition 5.1 in \\cite{Feireisl-2001}.\n\\begin{lemma}\\label{lem-viscous-flux}\nFor any $\\phi\\in \\mathcal{D}(0,T)$, $\\vp \\in \\mathcal{D}(\\O)$, we have that\n\\ba\\label{vis-flux}\\lim_{n\\to \\infty}\\int_0^T \\intO{\\phi(t)\\vp(x) \\left[p(\\vr_n) - \\left(\\frac{2\\mu^S(\\eta_n)}{3}-\\mu^B(\\eta_n)\\right) \\Div_x \\vu_n \\right] T_k(\\vr_n) }\\, \\dt\\\\\n=\\int_0^T \\intO{\\phi(t)\\vp(x) \\left[ \\overline{p(\\vr)} - \\left(\\frac{2\\mu^S(\\eta)}{3}-\\mu^B(\\eta)\\right) \\Div_x \\vu\\right] \\overline{T_k(\\vr)} }\\, \\dt,\n\\ea\nwhere $T_k$ is a cut-off function defined by $T_k(\\cdot) := k \\,T(\\frac{\\cdot}{k})$ for some concave function $T\\in C^\\infty([0,\\infty))$ such that $T(s)=s$ for $0\\leq s\\leq 1$ and $T(s)=2$ for $s\\geq 3$. Here $\\overline{T_k(\\vr)}$ is the weak* limit of the sequence $\\{T_k(\\vr_n)\\}_{n\\in \\N}$ in $L^\\infty((0,T)\\times \\O)$ as $n$ goes to infinity.\n\\end{lemma}\n\n\nThe viscosity coefficients $\\mu^B$ and $\\mu^S$ in our case are not constant, which gives rise to additional complications in the proof of this lemma. We shall employ the commutator estimates stated in Lemma \\ref{lem-Riesz3} and the techniques used in Section 3.6.5 in \\cite{F-N-book} to overcome the resulting difficulties.\n\n\n\\begin{proof}[Proof of Lemma \\ref{lem-viscous-flux}] As in the proof of the effective viscous flux lemma for the compressible Navier--Stokes equations with constant viscosity coefficients (see Proposition 5.1 in \\cite{Feireisl-2001} or Proposition 7.36 in \\cite{N-book}), we introduce the following test functions:\n\\be\\label{def-test-flux}\n\\vv_n(t,x) = \\phi(t) \\vp(x) \\A \\left[T_k(\\vr_n) \\right],\\quad \\vv(t,x) = \\phi(t) \\vp(x) \\A \\left[ \\overline{T_k(\\vr)}\\right],\n\\ee\nwhere $\\phi\\in C_c^\\infty(0,T), \\ \\vp \\in C_c^\\infty(\\O)$, and $\\A$ is the Fourier integral operator introduced in \\eqref{def-Reize}. We remark that $\\vr_n$ and $\\vr$ are extended by zero outside of $\\O$ in \\eqref{def-test-flux}.\n\n\nTaking these $\\vv_n$ and $\\vv$ as test functions in the weak formulation \\eqref{weak-form3} and the weak formulation of \\eqref{i2-0}, respectively, results in\n\\begin{align}\n\\label{weak-vn}\n&\\int_0^T \\intO{ \\big[ \\vr_n \\vu_n \\cdot \\partial_t \\vv_n + (\\vr_n \\vu_n \\otimes \\vu_n) : \\Grad \\vv_n + p(\\vr_n)\\, \\Div_x \\vv_n - \\SSS_n : \\Grad \\vv_n \\big]} \\,\\dt\n\\nonumber\\\\\n&\\hspace{3in}= \\int_0^T \\intO{ \\TT_n : \\nabla_x \\vv_n - \\vr_n\\, \\ff \\cdot \\vv_n} \\, \\dt,\\\\\n\\label{weak-v}\n&\\int_0^T \\intO{ \\left[ \\vr \\vu \\cdot \\partial_t \\vv + (\\vr \\vu \\otimes \\vu) : \\Grad \\vv + \\overline{p(\\vr)}\\, \\Div_x \\vv - \\SSS : \\Grad \\vv \\right]}\\,\\dt\n\\nonumber\\\\\n&\\hspace{3in}= \\int_0^T \\intO{ \\TT : \\nabla_x \\vv - \\vr\\, \\ff \\cdot \\vv} \\, \\dt.\n\\end{align}\n\nThe main idea is to pass to the limit $n\\to \\infty$ in \\eqref{weak-vn}, and compare the resulting limit with \\eqref{weak-v}, which will ultimately imply our desired result \\eqref{vis-flux}. Since, following the contributions of Lions \\cite{Lions-C} and Feireisl \\cite{F-book}, this type of argument is by now well understood, instead of including all of the details here we shall focus on the terms that do not appear in the case of the compressible Navier--Stokes equations with constant viscosity coefficients. These `new' terms are the following:\n\\ba\\label{flux-new-terms1}\n&\\int_0^T \\intO{ \\TT_n : \\nabla_x \\vv_n} \\, \\dt,\\quad \\int_0^T \\intO{ \\TT : \\nabla_x \\vv } \\, \\dt,\\quad \\int_0^T \\intO{ \\SSS_n : \\Grad \\vv_n } \\, \\dt,\\quad \\int_0^T \\intO{ \\SSS : \\Grad \\vv } \\, \\dt,\\nn\n\\ea\nand the goal is to prove that\n\\ba\\label{flux-new-terms2}\n\\lim_{n\\infty}\\int_0^T \\intO{ \\TT_n : \\nabla_x \\vv_n} \\, \\dt = \\int_0^T \\intO{ \\TT : \\nabla_x \\vv } \\, \\dt,\n\\ea\nand\n\\ba\\label{flux-new-terms3}\n&\\lim_{n\\to\\infty}\\int_0^T \\intO{ \\SSS_n : \\Grad \\vv_n } \\, \\dt - \\int_0^T \\intO{ \\SSS : \\Grad \\vv } \\, \\dt \\\\\n&= \\lim_{n\\to\\infty} \\int_0^T \\intO{\\phi(t)\\vp(x) \\left[ \\left(\\frac{2\\mu^S(\\eta_n)}{3}-\\mu^B(\\eta_n)\\right) \\Div_x \\vu_n \\right] T_k(\\vr_n) }\\, \\dt \\\\\n&\\quad - \\int_0^T \\intO{\\phi(t)\\vp(x) \\left[ \\left(\\frac{2\\mu^S(\\eta)}{3}-\\mu^B(\\eta)\\right) \\Div_x \\vu\\right] \\overline{T_k(\\vr)} }\\, \\dt.\n\\ea\n\nBy the strong convergence of $\\TT_n$ stated in \\eqref{TTn-lim-strff} in Section \\ref{sec:lim-TT} it is straightforward to show the convergence result \\eqref{flux-new-terms2}, so we only focus on proving \\eqref{flux-new-terms3}. We begin by noting that\n\\ba\\label{flux-new31}\n\\int_0^T \\intO{ \\SSS_n : \\Grad \\vv_n } \\, \\dt = \\int_0^T \\phi(t)\\intO{ \\SSS_n :\\, \\left(\\nabla_x \\vp \\otimes \\A[ T_k(\\vr_n)]\\right)} \\, \\dt + \\int_0^T \\phi(t)\\intO{ \\vp \\, \\SSS_n : \\RR [ T_k(\\vr_n)]} \\, \\dt,\\nn\n\\ea\nwhere $\\RR$ is the Riesz operator defined in \\eqref{def-Reize} in Section \\ref{sec:Riesz}.\n\n\nFor any fixed $k\\in\\N$, the sequence $\\{ T_k(\\vr_n)\\}_{n\\in \\N}$ is uniformly bounded in $ L^\\infty(0,T;L^r(\\R^3))$ for all $r \\in [1, \\infty]$. Thus, by Lemma \\ref{lem-Riesz1}, we have that\n\\be\\label{A-Tk-vrn-1}\n\\sup_{n\\in \\N} \\left\\| \\A[T_k(\\vr_n)]\\right\\|_{L^\\infty(0,T;W^{1,r}(\\O;\\R^3))} \\leq c\\quad \\mbox{for any $r \\in (1,\\infty)$}.\n\\ee\n\nObserving that $T_k(\\cdot)$ fulfills the properties in \\eqref{cond-b-renormal}, it follows from \\eqref{weak-renormal1} that\n\\be\\label{renomal-Tk}\n\\d_t T_k(\\vr_n) + \\Div_x\\left(T_k(\\vr_n)\\vu_n\\right) + \\left(T_k'(\\vr_n)\\vr_n-T_k(\\vr_n)\\right)\\Div_x \\vu_n =0\\quad \\mbox{in $ \\mathcal{D}'((0,T)\\times \\R^3)$}.\n\\ee\nThis implies that\n\\ba\\label{A-Tk-vrn-2}\n\\d_t \\A[ T_k(\\vr_n)]= \\A[\\d_t T_k(\\vr_n)] = - \\A\\left[ \\Div_x\\left(T_k(\\vr_n)\\vu_n\\right) \\right] - \\A\\left[ \\left(T_k'(\\vr_n)\\vr_n-T_k(\\vr_n)\\right)\\Div_x \\vu_n \\right].\n\\ea\nSince\n$\n\\A_j \\Div_x = -\\sum_{i=1}^3\\d_j \\d_i \\Delta^{-1} = - \\sum_{i=1}^3\\RR_{ij}\n$\nare Riesz type operators, by Lemma \\ref{lem-Riesz1} we have that\n$$\n\\|\\A_j \\left[ \\Div_x \\vvarphi \\right]\\|_{L^r(\\R^3:\\R^3)}\\leq c \\, \\|\\vvarphi\\|_{L^r(\\R^3;\\R^3)} \\quad \\mbox{for any $\\vvarphi\\in L^r(\\R^3;\\R^3)$ and any $r \\in (1,\\infty)$}.\n$$\nBy using \\eqref{vun-limit1} and \\eqref{vun-limit2} we deduce that\n\\ba\\label{A-Tk-vrn-3}\n\\left\\|\\A\\left[ \\Div_x\\left(T_k(\\vr_n)\\vu_n\\right) \\right]\\right\\|_{L^2(0,T;L^6(\\O;\\R^3))} &\\leq c &&\\quad \\mbox{when $\\o\\geq 0$}; \\\\\n\\left\\|\\A\\left[ \\Div_x\\left(T_k(\\vr_n)\\vu_n\\right) \\right]\\right\\|_{L^2(0,T;L^{\\frac{12}{2+6|\\o|}}(\\O;\\R^3))} &\\leq c &&\\quad\\mbox{when $\\o\\leq 0$}.\n\\ea\nFurthermore, by Lemma \\ref{lem-Riesz1} in conjunction with \\eqref{vun-limit1} and \\eqref{vun-limit2} we deduce that\n\\ba\\label{A-Tk-vrn-4}\n\\left\\| \\A\\left[ \\left(T_k'(\\vr_n)\\vr_n-T_k(\\vr_n)\\right)\\Div_x \\vu \\right] \\right\\|_{L^2(0,T;L^6(\\O;\\R^3))} &\\leq c &&\\quad \\mbox{when $\\o\\geq 0$}; \\\\\n\\left\\| \\A\\left[ \\left(T_k'(\\vr_n)\\vr_n-T_k(\\vr_n)\\right)\\Div_x \\vu \\right] \\right\\|_{L^2(0,T;L^{\\frac{12}{2+6|\\o|}}(\\O;\\R^3))} &\\leq c &&\\quad \\mbox{when $\\o\\leq 0$}.\n\\ea\nThus,\n\\ba\\label{A-Tk-vrn-5}\n\\left\\| \\d_t \\A[ T_k(\\vr_n)] \\right\\|_{L^2(0,T;L^6(\\O;\\R^3))} &\\leq c &&\\quad \\mbox{when $\\o\\geq 0$};\\\\\n\\left\\| \\d_t \\A[ T_k(\\vr_n)] \\right\\|_{L^2(0,T;L^{\\frac{12}{2+6|\\o|}}(\\O;\\R^3))} &\\leq c &&\\quad \\mbox{when $\\o\\leq 0$}.\n\\ea\nIt follows from the uniform estimates in \\eqref{A-Tk-vrn-1} and \\eqref{A-Tk-vrn-5} and the Aubin--Lions--Simon compactness theorem that\n\\ba\\label{A-Tk-vrn-f}\n\\A[T_k(\\vr_n)] \\to \\overline{\\A[T_k(\\vr)]} = \\A[\\overline{T_k(\\vr)}] \\quad \\mbox{strongly in $L^\\infty(0,T;L^r(\\O;\\R^3))$ \\quad for any $r\\in (1,\\infty)$}.\n\\ea\nTogether with the weak convergence of $\\SSS_n$ stated in \\eqref{Newtonian-lim1} and \\eqref{Newtonian-lim10}, we deduce that\n\\ba\\label{flux-new32}\n \\lim_{n\\to\\infty}\\int_0^T \\phi(t)\\intO{ \\SSS_n : (\\nabla_x \\vp \\otimes \\A[ T_k(\\vr_n)])} \\, \\dt = \\int_0^T \\phi(t)\\intO{ \\SSS : (\\nabla_x \\vp \\otimes \\A[ \\overline{T_k(\\vr)}])} \\, \\dt.\n\\ea\nTherefore, to show \\eqref{flux-new-terms3}, it suffices to prove that\n\\ba\\label{flux-new33}\n& \\lim_{n\\to \\infty} \\int_0^T \\phi(t)\\intO{ \\vp \\, \\SSS_n : \\RR [ T_k(\\vr_n)]} \\, \\dt - \\int_0^T \\phi(t)\\intO{ \\vp \\, \\SSS : \\RR [ \\overline{T_k(\\vr)}]} \\, \\dt\\\\\n &= \\lim_{n\\to\\infty} \\int_0^T \\intO{\\phi(t)\\vp(x) \\left[ \\left(\\frac{2\\mu^S(\\eta_n)}{3}-\\mu^B(\\eta_n)\\right) \\Div_x \\vu_n \\right] T_k(\\vr_n) }\\, \\dt \\\\\n&\\quad - \\int_0^T \\intO{\\phi(t)\\vp(x) \\left[ \\left(\\frac{2\\mu^S(\\eta)}{3}-\\mu^B(\\eta)\\right) \\Div_x \\vu\\right] \\overline{T_k(\\vr)} }\\, \\dt.\n\\ea\n\nBy \\eqref{i3} we have that\n\\ba\\label{flux-new34}\n&\\int_0^T \\phi(t)\\intO{ \\vp \\, \\SSS_n : \\RR [ T_k(\\vr_n)]} \\ \\dt= \\int_0^T \\phi(t)\\intO{ \\vp \\, \\mu^S(\\eta_n) \\left( \\frac{\\nabla \\vu_n + \\nabla^{\\rm T} \\vu_n}{2} \\right) : \\RR [ T_k(\\vr_n)]} \\, \\dt\\\\\n&\\qquad + \\int_0^T \\phi(t)\\intO{ \\vp \\, \\left(\\mu^B(\\eta_n)-\\frac{\\mu^S(\\eta_n)}{3}\\right) (\\Div_x \\vu_n) \\II: \\RR [ T_k(\\vr_n)]} \\, \\dt\\\\\n& = \\int_0^T \\phi(t)\\intO{ \\vp \\, \\mu^S(\\eta_n) \\nabla \\vu_n : \\RR [ T_k(\\vr_n)]} \\, \\dt + \\int_0^T \\phi(t)\\intO{ \\vp \\, \\left(\\mu^B(\\eta_n)-\\frac{\\mu^S(\\eta_n)}{3}\\right) (\\Div_x \\vu_n) \\, [ T_k(\\vr_n)]} \\, \\dt,\n\\ea\nwhere have we used the fact that $\\RR$ is symmetric (see Lemma \\ref{lem-Riesz1}) and that $\\II : \\RR = \\sum_{i=1}^3 \\RR_{ii}=\\II$.\nFurther, by Lemma \\ref{lem-Riesz1}, and noting that $\\sum_{i,j=1}^3 \\RR_{ij}\\partial_j \\vu_n^i = \\Div_x \\vu_n$, we have that\n\\ba\\label{flux-new35}\n& \\int_0^T \\phi(t)\\intO{ \\vp \\, \\mu^S(\\eta_n) \\nabla \\vu_n : \\RR [ T_k(\\vr_n)]} \\, \\dt = \\int_0^T \\phi(t)\\intO{ \\sum_{i,j=1}^3 \\RR_{ij} \\left[\\vp \\, \\mu^S(\\eta_n)\\, \\d_j \\vu_n^i \\right] T_k(\\vr_n)} \\, \\dt\\\\\n&\\quad =\\int_0^T \\phi(t)\\intO{ \\sum_{i,j=1}^3 \\left( \\vp \\, \\mu^S(\\eta_n)\\, \\RR_{ij} \\left[ \\d_j \\vu_n^i \\right] T_k(\\vr_n) + R_n^{ij} \\, T_k(\\vr_n) \\right)\\!} \\, \\dt \\\\\n&\\quad = \\int_0^T \\phi(t)\\intO{ \\left(\\vp \\, \\mu^S(\\eta_n)\\, (\\Div_x \\vu_n) \\, T_k(\\vr_n) + \\sum_{i,j=1}^3 R_n^{ij} \\, T_k(\\vr_n) \\right)\\!} \\, \\dt,\n\\ea\nwhere $R_n^{ij}$ is the commutator defined by $R_n^{ij}:= \\RR_{ij} \\left[\\vp \\, \\mu^S(\\eta_n) \\d_j \\vu_n^i \\right]- \\vp \\, \\mu^S(\\eta_n) \\RR_{ij} \\left[ \\d_j \\vu_n^i \\right].$ We will use Lemma \\ref{lem-Riesz3} to derive uniform bounds on $R_n^{ij}$. We first consider the case $0\\leq \\o < \\frac 53$. Then, by \\eqref{mu-eta1} and \\eqref{etan-est1}, the sequence\n$\\{\\nabla_x \\mu^S(\\eta_n)\\}_{n\\in\\N} = \\{(\\mu^S)'(\\eta_n) \\nabla_x \\eta_n\\}_{n\\in\\N}$ is uniformly bounded in $L^2(0,T;L^2(\\O;\\R^3))$ when $0\\leq \\o\\leq 1$, and in $L^{2}(0,T;L^\\frac{2}{\\o}(\\O;\\R^3))$ when $1\\leq \\o < \\frac 53$.\nThus,\n$$\n\\sup_{n\\in \\N}\\| \\mu^S(\\eta_n) \\|_{L^2(0,T;W^{1,2}(\\O))} \\leq c \\quad \\mbox{when}\\ 0\\leq \\o\\leq 1;\\qquad \\sup_{n\\in \\N}\\| \\mu^S(\\eta_n) \\|_{L^2(0,T;W^{1,\\frac{2}{\\o}}(\\O))} \\leq c \\quad \\mbox{when}\\ 1 \\leq \\o < \\frac 53.\n$$\n\nBy \\eqref{vun-limit1} the sequence $\\{\\nabla_x \\vu_n\\}_{n\\in\\N}$ is uniformly bounded in $ L^2(0,T; L^{2}(\\O;\\R^{3 \\times 3}))$. Furthermore, we note that\n\\ba\\label{mu-eta10}\n\\frac{1}{2} + \\frac{1}{2} -\\frac{1}{3} = \\frac{2}{3}< 1;\\quad \\frac{1}{2}+\\frac{\\o}{2} - \\frac{1}{3} =\\frac{1+3\\o}{6} < 1, \\quad \\mbox{as long as}\\ 0\\leq \\o < \\frac{5}{3}.\n\\ea\nHence, by Lemma \\ref{lem-Riesz3}, for any $s>1$ such that\n\\ba\\label{est-mun-2}\n \\frac{1}{s} >\\frac{1}{2}+\\frac{1}{2} - \\frac{1}{3}=\\frac{2}{3} \\quad \\mbox{when}\\ 0\\leq \\o\\leq 1;\\quad \\frac{1}{s} >\\frac{1}{2}+\\frac{\\o}{2} - \\frac{1}{3} =\\frac{1+3\\o}{6} \\quad \\mbox{when}\\ 1\\leq \\o < \\frac{5}{3},\n\\ea\nwe have, for $r_\\o=2$ when $0\\leq \\o\\leq 1$ and for $r_\\o=\\frac{2}{\\o}$ when $1\\leq \\o < \\frac{5}{3}$, that\n\\ba\\label{flux-commu-2}\n&\\left\\|R_n^{ij}\\right\\|_{W^{\\b,s}(\\R^3)} \\leq c\\, \\| \\vp \\mu^S(\\eta_n) \\|_{W^{1,r_\\o}(\\O)} \\|\\d_j \\vu_n^i\\|_{L^2(\\O)},\n\\ea\nwhere $\\b\\in (0,1)$ is such that\n\\ba\\label{est-mun-3}\n\\frac{\\b}{3} =\\frac 13 + \\frac 1s - 1 \\quad \\mbox{when}\\ 0\\leq \\o\\leq 1;\\quad \\frac{\\b}{3} =\\frac 13 + \\frac 1s - \\frac{1+\\o}{2}\\quad \\mbox{when}\\ 1\\leq \\o < \\frac 53.\n\\ea\nTherefore, for some $s>1$ and some $0<\\b<1$, we have the uniform estimate\n\\ba\\label{flux-commu-3}\n\\sup_{n\\in \\N} \\left\\|R_n^{ij} \\right\\|_{L^1(0,T;W^{\\b,s}(\\R^3))} \\leq c\\, \\sup_{n\\in \\N} \\| \\vp \\mu^S(\\eta_n) \\|_{L^2(0,T;W^{1,r_\\o}(\\O))}\\, \\sup_{n\\in \\N} \\|\\d_j \\vu_n^i\\|_{L^2(0,T;L^2(\\O))}\\leq c.\n\\ea\nHence, by the strong convergence of $\\eta_n$ shown in Section \\ref{sec:lim-eta} we have that\n\\ba\\label{flux-commu-4}\nR_n^{ij} \\to R^{ij}:= \\RR_{ij} \\left[\\vp \\, \\mu^S(\\eta) \\d_j \\vu^i \\right]- \\vp \\, \\mu^S(\\eta) \\RR_{ij} \\left[ \\d_j \\vu^i \\right] \\quad \\mbox{weakly in} \\ L^{\\frac{10}{3\\o+5}}((0,T)\\times \\O).\n\\ea\n\nBy the boundedness of $T_k$ we have\n\\be\\label{Tk-weak}\nT_k(\\vr_n)\\to \\overline{T_k(\\vr)} \\quad \\mbox{weakly in $L^r((0,T)\\times \\O)$\\quad for any $11$ such that\n\\ba\\label{est-mun-2-w<0}\n& \\frac{1}{s} > \\frac{1}{2}+\\frac{2+|\\o|}{4} - \\frac{1}{3},\n\\ea\nwe have that\n\\ba\\label{flux-commu-2-w<0}\n&\\left\\|R_n^{ij}\\right\\|_{W^{\\b,s}(\\R^3)} \\leq c\\, \\| \\vp \\mu^S(\\eta_n) \\|_{W^{1,2}(\\O)} \\|\\d_j \\vu_n^i\\|_{L^\\frac{4}{2+|\\o|}(\\O)},\n\\ea\nwhere $\\b\\in (0,1)$ is such that\n\\ba\\label{est-mun-3-w<0}\n\\frac{\\b}{3} =\\frac 13 + \\frac 1s - \\frac{1}{2}-\\frac{2+|\\o|}{4}.\n\\ea\nTherefore, for some $s>1$ and some $\\b>0$, the following uniform estimate holds:\n\\ba\\label{flux-commu-3-w<0}\n\\sup_{n\\in \\N} \\left\\|R_n^{ij} \\right\\|_{L^1(0,T;W^{\\b,s}(\\R^3))} \\leq c\\, \\sup_{n\\in \\N} \\| \\vp \\mu^S(\\eta_n) \\|_{L^2(0,T;W^{1,2}(\\O))}\\,\\sup_{n\\in \\N} \\|\\d_j \\vu_n^i\\|_{L^2(0,T;L^\\frac{4}{2+|\\o|}(\\O))}\\leq c.\n\\ea\nThanks to the strong convergence of $\\eta_n$ shown in Section \\ref{sec:lim-eta}, we have that\n\\ba\\label{flux-commu-4-w<0}\n R_n^{ij} \\to R^{ij}:= \\RR_{ij} \\left[\\vp \\, \\mu^S(\\eta) \\d_j \\vu^i \\right]- \\vp \\, \\mu^S(\\eta) \\RR_{ij} \\left[ \\d_j \\vu^i \\right], \\\\\n \\mbox{weakly in} \\ L^{r}((0,T)\\times \\O),\\ \\mbox{for any $r<\\frac{20}{10+3|\\o|}$}.\n\\ea\nThus, the Div-Curl lemma implies that\n\\ba\\label{flux-commu-5-w<0}\nR_n^{ij}\\, T_k(\\vr_n) \\to R^{ij} \\overline{T_k(\\vr_n)} \\quad \\mbox{weakly in} \\ L^{r}(0,T;L^{r}(\\O)) \\quad \\mbox{for any $r<\\frac{20}{10+3|\\o|}$},\n\\ea\nwhich, once again, implies \\eqref{flux-new36}.\n\n\\medskip\n\n\n\n\nFinally, from \\eqref{flux-new34}, \\eqref{flux-new35} and \\eqref{flux-new36}, we deduce the desired result \\eqref{flux-new33}. This implies \\eqref{flux-new-terms3} by using \\eqref{flux-new32}. At the same time, by tedious but, by now, standard calculations, as in the proof of the effective viscous flux lemma in \\cite{Lions-C, F-book, N-book}, we have that\n\\ba\\label{flux-new-terms3f}\n\\lim_{n\\to\\infty}\\int_0^T \\intO{ \\SSS_n : \\Grad \\vv_n } \\ \\dt - \\int_0^T \\intO{ \\SSS : \\Grad \\vv } \\ \\dt =0.\n\\ea\nCombining \\eqref{flux-new-terms3f} with \\eqref{flux-new-terms3}, we obtain \\eqref{vis-flux} and complete the proof of Lemma \\ref{lem-viscous-flux}.\n\\end{proof}\n\n\n\\medskip\n\nWith Lemma \\ref{lem-viscous-flux} in hand, the proof of the strong convergence result \\eqref{pn-limf2} then proceeds along a well-understood route (see for example \\cite{Feireisl-2001}, from Section 6 to Section 8), so we shall not dwell on the details here. In particular, as in Proposition 7.1 in \\cite{Feireisl-2001}, the limit $(\\vr, \\vu)$ satisfies \\eqref{weak-renormal1} in the sense of renormalized weak solutions.\n\n\\medskip\n\nIt remains to show that the limit $(\\vr,\\vu,\\psi,\\eta)$ satisfies the energy inequality \\eqref{energy}. This is easily seen by noting the strong convergence assumption on the initial data in \\eqref{cond-ini-n} and passage to the limit $n\\to \\infty$ in \\eqref{energy-n}, which directly imply the energy\ninequality \\eqref{energy} \nby the application of Tonelli's sequential weak (weak*) lower semicontinuity theorem (cf. Theorem 10.16 in \\cite{RR-book}, for example,) to the terms appearing on the left-hand side of \\eqref{energy-n}; in particular, the sequential weak lower semicontinuity of the $L^p$ norm, $1< p<\\infty$, the sequential weak* lower semicontinuty of the $L^\\infty$ norm and inequality \\eqref{psin-est3} are used.\nThus we have shown that the limit $(\\vr,\\vu,\\psi,\\eta)$ is a dissipative (finite-energy) weak solution in the sense of Definition \\ref{def-weaksl}. That completes the proof of Theorem \\ref{theorem}.\n\n\\begin{remark}\\label{rem-mu-eta}\nThe constraint $\\o<\\frac{5}{3}$ for the case $\\o\\geq 0$ is crucially determined by \\eqref{mu-eta10} and the condition $\\frac{10}{3\\o+5}>1$, which appears in \\eqref{flux-commu-4}, while the constraint $\\o>-\\frac{4}{3}$ for the case $\\o\\leq 0$ is crucially determined by \\eqref{mu-eta10-w<0}. It is unclear whether,\nwith our present techniques at least, these restrictions on $\\o$ can be relaxed.\n\\end{remark}\n\n\n\n\\section{Conclusion: existence of dissipative weak solutions}\\label{sec:end}\n\n\\label{End}\n\nThe conclusions of Theorem \\ref{theorem} \\emph{do not}, of course, imply the existence of dissipative (finite-energy) weak solutions to the Navier--Stokes--Fokker--Planck system with polymer-number-density-dependent viscosity coefficients. A rigorous proof of the existence of dissipative weak solutions would require the following:\n\\begin{itemize}\n\n\\item a suitable approximation scheme, compatible with the energy inequality and the compactness arguments presented in this paper;\n\n\\item a rigorous proof of the existence of a solution to the approximation scheme;\n\n\\item proof of the convergence of the sequence of approximate solutions to a dissipative weak solution, mimicking the arguments presented in this paper.\n\n\\end{itemize}\n\nGiven the formal similarity of the present model to the one studied in \\cite{Barrett-Suli}, a natural approach would be to adjust the approximation scheme used\nin \\cite{Barrett-Suli}, based on time-discretization, in the case of the Navier--Stokes--Fokker--Planck system with constant viscosity coefficients (or, alternatively, to use a Galerkin approximation scheme in the spatial variables, similar to the one in \\cite{FNP}). The added technical difficulties, caused by the presence of the variable viscosity coefficients, can be handled exactly as in \\cite{AbFei}, where a similar scheme, based on time-discretization, was applied to a diffuse interface model with viscosity coefficients that depended on the concentration.\n\n\n\n\\section*{Acknowledgements}\n\\thispagestyle{empty}\n\nThe authors acknowledge the support of the project LL1202 in the programme ERC-CZ funded by the Ministry of Education, Youth and Sports of the Czech Republic.\n\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nThe discovery of the standard model (SM) Higgs boson~\\cite{Chatrchyan:2012xdj,Aad:2012gk}, \nthe only known elementary scalar particle, has been the main result of the Large Hadron Collider\nphysics program. Remarkably, the measured Higgs boson mass was actually predicted with a good accuracy already in 1995~\\cite{Froggatt:1995rt} by\napplying the principle of multiple point criticality (PMPC)~\\cite{Bennett:1993pj,Bennett:1996vy,Bennett:1996hx} to the Higgs boson effective potential. \nToday the SM Higgs criticality has been studied to three-loop accuracy, indicating that the present data prefers a meta-stable Higgs vacuum \nsufficiently close to criticality~\\cite{Buttazzo:2013uya}.\\footnote{This statement depends most strongly on the precise value of the top quark Yukawa coupling, as the critical point remains within $3\\sigma$ uncertainty of $y_t$. Since $y_t$ is not a directly measured quantity, the systematic uncertainty related to its derivation from collider top mass measurements may be presently underestimated~\\cite{Corcella:2019tgt}.}\nThis fact is the most interesting unexplained experimental outcome of collider physics so far.\n\n\nWithin the SM, there is no explanation as to why the SM parameters -- the gauge, Yukawa and scalar quartic couplings -- must take values leading to the Higgs criticality.\nTherefore, in the original proposal~\\cite{Froggatt:1995rt}, the existence of multiple degenerate vacua was argued to be analogous to the \nsimultaneous existence of multiple phases in thermodynamics. This statement is purely empirical and cannot follow solely from the SM physics, nevertheless\nit may have some support from statistical systems in solid state physics~\\cite{Volovik:2003hn}.\nTo derive the PMPC from fundamental physical principles, Higgs criticality has been attributed to the proposed asymptotic safety of quantum gravity~\\cite{Shaposhnikov:2009pv}.\nMore recently Dvali~\\cite{Dvali:2011wk} revived the old argument by Zeldovich~\\cite{Zeldovich:1974py,Kobzarev:1974cp} that physical systems with multiple vacua must all be degenerate because of the consistency of quantum field theory -- tunneling from a Minkowski vacuum to a lower one must be exponentially fast as in any theory possessing ghosts. \nThe Zeldovich-Dvali argument contradicts Coleman's semi-classical computation in the Euclidian space that the formation of vacuum bubbles is exponentially slow~\\cite{Coleman:1977py}, and has also been challenged by studies of different physical systems~\\cite{Garriga:2011we}. However, attempts to generalize Coleman's computation to the realistic case with zero-mass vacuum bubbles, which is the origin of the divergence, have failed and the Zeldovich-Dvali argument has not been disproved~\\cite{Gross:2020tph}. Altogether, the underlying fundamental physics behind the PMPC remains presently unclear: we accept it on empirical grounds.\n\nNevertheless, the phenomenological success of PMPC has motivated different studies of scalar potentials in extended models. Already the very first papers\nafter the Higgs boson discovery argued that if the SM is extended by a singlet scalar, the Higgs vacuum stability as well as criticality can be achieved due to the radiative corrections of the singlet~\\cite{Kadastik:2011aa}. Later the PMPC has been applied to the singlet~\\cite{Haba:2014sia,Kawana:2014zxa,Haba:2016gqx,McDowall:2019knq} and to the two Higgs doublet models (2HDM)~\\cite{Froggatt:2004st,Froggatt:2006zc,Froggatt:2008am,McDowall:2018ulq,Maniatis:2020wfz}.\nThe Higgs criticality may play an important role in understanding Higgs inflation~\\cite{Bezrukov:2007ep}. We do not make an attempt to form a comprehensive list of all possible applications of PMPC and limit ourselves to the above examples.\n\nThe aim of this work is to point out that so far all the applications of PMPC to scalar potentials, including the examples given above, deal with radiatively generated\nmultiple minima. However, if the degeneracy of multiple vacua is a fundamental property of Nature, all vacua of multi-scalar models -- including the ones \ndominated by {\\it tree-level} potential terms -- should be degenerate. We first argue that this formulation of PMPC should be correct, and after that exemplify its physics potential, analyzing the vacuum solutions of the scalar singlet dark matter (DM) model~\\cite{Silveira:1985rk,McDonald:1993ex,Burgess:2000yq,Cline:2013gha} (see~\\cite{Arcadi:2019lka} for a review).\nWe present three simple scalar DM models utilizing the PMPC: I) $\\mathbb{Z}_2$-symmetric real scalar DM model, II) $\\mathbb{Z}_2$-symmetric pseudo-Goldstone DM (pGDM) model and III) $\\mathbb{Z}_3$-symmetric pseudo-Goldstone DM model.\\footnote{In a more complicated pGDM model, all these global symmetries can be broken, save for the CP-like $S \\to S^{*}$ that stabilizes DM \\cite{Alanne:2020jwx}.} Our results will show that the allowed parameter spaces of those models are constrained and the PMPC \nwill, in principle, be testable if the DM of Universe is the singlet scalar. Generalizing this result, our work demonstrates that the PMPC might be a useful tool to discriminate\nbetween different models of new physics beyond the SM.\n\nThis work is organized as follows. In Section~\\ref{sec:PMPC} we generalize the concept of PMPC to the case of multiple scalar fields. In Sections~\\ref{sec:model1},~\\ref{sec:model2} and~\\ref{sec:model3} we analyze our example models. We conclude and discuss perspectives in Section~\\ref{sec:conc}.\n\n\n\n\\section{Tree-level PMPC}\n\\label{sec:PMPC}\n\nSo far the PMPC has been used to constrain model parameters by requiring that radiatively generated minimum of the\nscalar potential is degenerate with the electroweak vacuum. This was the case with the SM as well as with the scalar singlet and doublet\nextensions of the SM. Clearly, in multi-scalar models there can be several minima of the \nscalar potential already at tree level. If the PMPC follows from some underlying fundamental physics principle, all those vacua must be degenerate as well.\nThis is the concept we adopt and study in this work. Although we do not know which physics principle is behind the PMPC, we exemplify with the following discussion\nthat such principles may exist.\n\nAs Zeldovich pointed out in early 70ies~\\cite{Zeldovich:1974py,Kobzarev:1974cp}, if there exist two minima of the scalar potential with $V=0$ and $V<0$,\nthe vacuum decay bubble with mass $m = 0$ can appear with any initial velocity, giving rise to a divergent boost integral. It should be noticed that \nthere is no such a problem with the pair creation of ordinary particles. Instead, this problem is intrinsic to the theories with ghosts, {\\it i.e.}, to particles with negative energy. \nWhile the theories with ghosts are widely believed to be pathological precisely because of the exponentially fast vacuum decay, similar phenomenon\nshould appear in the Lorentz-invariant decay of the Minkowski vacuum, rendering this configuration unphysical. Dvali, who apparently was not aware of the \nZeldovich paper, independently make this claim in Ref.~\\cite{Dvali:2011wk}. More recently this setup was studied in theories of ghosts~\\cite{Gross:2020tph}, finding the same divergence.\n\n\nNeedless to say, this claim contradicts Coleman's result that the vacuum decay rate is finite and exponentially suppressed~\\cite{Coleman:1977py}. \nAs Coleman computed one bubble decay rate in Euclidean space, the authors of Ref.~\\cite{Gross:2020tph} wanted to generalize Coleman's computation to\nthe case of ghosts, but did not succeed. The Zeldovich-Dvali claim was also criticized in Ref.~\\cite{Garriga:2011we} which,\nhowever, studies a completely different physical situation -- particle pair production in uniform electric field. Whether this set-up is analogous to the case of \nvacuum decay is far from being obvious. It seems that the main source of different opinions is whether the divergent boost integral has to be taken or not,\nand there is no definitive answer at the moment. \n\nNevertheless, this discussion demonstrates that there exist attempts to derive PMPC from general physics principles.\nThe PMPC in this, broader formulation has been used to criticize certain supersymmetry breaking scenarios~\\cite{Bajc:2011iu}.\nIn this work we perform phenomenological analyses of implications of tree-level PMPC on multi-scalar DM models.\n\n\n\\section{Model I: $\\mathbb{Z}_2$ real singlet}\n\\label{sec:model1}\n\nThe real scalar singlet DM model extends the SM with one real scalar singlet which is stabilized by a $\\mathbb{Z}_2$ symmetry. \nThis is among the most popular DM models whose phenomenology has been thoroughly studied~\\cite{Davoudiasl:2004be,Ham:2004cf,OConnell:2006rsp,Patt:2006fw,Profumo:2007wc,Barger:2007im,He:2007tt,He:2008qm,Yaguna:2008hd,Lerner:2009xg,Farina:2009ez,Goudelis:2009zz,Profumo:2010kp,Guo:2010hq,Barger:2010mc,Arina:2010rb,Bandyopadhyay:2010cc,Drozd:2011aa,Djouadi:2011aa,Low:2011kp,Mambrini:2011ik,Espinosa:2011ax,Mambrini:2012ue,Djouadi:2012zc,Cheung:2012xb,Cline:2013gha,Urbano:2014hda,Endo:2014cca,Feng:2014vea,Duerr:2015bea,Duerr:2015mva,Beniwal:2015sdl,Cuoco:2016jqt,Escudero:2016gzx,Han:2016gyy,He:2016mls,Ko:2016xwd,Athron:2017kgt,Ghorbani:2018yfr}. Nevertheless, we shall show that applying the PMPC to this model scalar potential will open new perspectives in\nunderstanding the allowed parameter space of the model.\n\nThe scalar potential of the model is given by\n\\begin{equation}\nV=\\mu_h^2 \\lvert H \\rvert^{2} + \\frac{\\mu_s^2}{2}s^2+\\lambda_h \\lvert H \\rvert^{4}+\\frac{\\lambda_s}{4}s^4+\\frac{\\lambda_{hs}}{2} \\lvert H \\rvert^{2} s^2,\n\\end{equation}\nwhere $H$ is the SM Higgs doublet in unitary gauge:\n\\begin{equation}\nH=\\frac{1}{\\sqrt{2}}\\left(\\begin{array}{c}\n0\\\\\nh+v_h\n\\end{array}\n\\right).\n\\end{equation}\nThe Higgs vacuum expectation value is $v_h=246$ GeV.\nThe scalar $s$, as the dark matter candidate, must not acquire a non-zero VEV: then the scalar $s$ does not mix with the Higgs and remains stable due to the $\\mathbb{Z}_2$ symmetry. The freeze-out of the $s$ is independent of the self-coupling of the scalar singlet $\\lambda_s$ and its relic abundance is therefore set only by the Higgs portal coupling $\\lambda_{hs}$. \n\nThe potential may have three types of minima in the field space $(h,s)$: the Higgs vacuum $(v_h,0)$, singlet vacuum $(0,v_s)$ and a general or mixed vacuum $(v'_h,v'_s)$. Presently, the Universe is in the Higgs vacuum. According to the PMPC these minima are to be degenerate, if they exist simultaneously. Let us examine the three vacuum configurations more closely. \n\n\\subsection{The Higgs vacuum configuration $(v_h,0)$}\nThe Higgs vacuum $(v_h,0)$, in which we live, satisfies the following minimization condition:\n\\begin{equation}\n\\mu_h^2=-\\lambda_h v_h^2.\\label{Our vacuum minimization condition}\n\\end{equation} \nThe scalar mass matrix is given by\n\\begin{equation}\n\\mathcal{M}^2=\\left(\\begin{array}{cc}\n2\\lambda_h v_h^2 & 0\\\\\n0 & \\frac{1}{2}\\lambda_{hs} v_h^2+\\mu_s^2\n\\end{array}\n\\right).\n\\end{equation}\n\n\n\\subsection{The singlet vacuum configuration $(0,v_s)$}\nThe minimization condition for the singlet vacuum $(0,v_s)$ is given by\n\\begin{equation}\n\\mu_s^2=-\\lambda_s v_s^2.\\label{singlet vacuum minimization condition}\n\\end{equation}\n\nIn this case, the scalar mass matrix is given by\n\\begin{equation}\n\\mathcal{M}^2=\\left(\\begin{array}{cc}\n\\frac{1}{2}\\lambda_{hs} v_s^2+\\mu_h^2 & 0\\\\\n0 & 2\\lambda_s v_s^2\n\\end{array}\n\\right).\n\\end{equation}\n\n\n\n\\subsection{The mixed vacuum configuration $(v'_h,v'_s)$}\nThe minimization conditions for the mixed vacuum $(v'_h,v'_s)$ are given by\n\\begin{eqnarray}\n\\mu_h^2 & = & -\\lambda_h {v'_h}^2-\\frac{1}{2}\\lambda_{hs}{v'_s}^2\\\\\n\\mu_s^2 & = & -\\lambda_{s} {v'_s}^2-\\frac{1}{2}\\lambda_{hs} {v'_h}^2\n\\end{eqnarray}\nThe scalar mass matrix is given by\n\\begin{equation}\n\\mathcal{M}^2=\\left(\\begin{array}{cc}\n2\\lambda_{h} {v'_h}^2 & \\lambda_{hs} {v'_h}v'_s\\\\\n\\lambda_{hs} {v'_h}v'_s & 2\\lambda_{s} {v'_s}^2\n\\end{array}\n\\right),\n\\end{equation}\nleading to mixing of the Higgs doublet with the real singlet.\n\n\\subsection{Comparison of minima}\nThe current minimum configuration is the Higgs vacuum $(v_h,0)$, which allows dark matter and correctly breaks the electroweak symmetry. If the other vacuum configurations are to exist simultaneously with the $(v_h,0)$, they must be degenerate. Imposing that the $(0,v_s)$ vacuum configuration is degenerate to $(v_h,0)$,\n\\begin{equation}\nV(h=v_h,s=0)= V(h=0,s=v_s),\n\\end{equation}\nleads to the condition\n\\begin{equation}\\label{the constraint for our vacuum VS singlet vacuum}\nm_s^4-\\lambda_{hs} m_s^2 v_h^2 -\\lambda_s\\lambda_h v_h^4+\\frac{1}{4}\\lambda_{hs}^2 v_h^4=0.\n\\end{equation}\n\nDemanding that the $(v'_h,v'_s)$ vacuum configuration be degenerate to the present Higgs vacuum $(v_h,0)$ leads to a contradiction. Therefore, according to the PMPC, there can at most only be two degenerate vacua: the Higgs vacuum $(v_h,0)$ and the singlet vacuum $(0,v_s)$. \n\nThe collider experiments impose constraints on the singlet $s$ when it is lighter than $m_h\/2$. Then the Higgs decay width to the $ss$ final state is \\cite{Cline:2013gha}\n\\begin{equation}\n\\Gamma (h\\to ss) =\\frac{\\sqrt{m_h^2-4m_s^2}}{32\\pi m_h^2}\\lambda_{hs}^2 v_h^2. \n\\end{equation}\nIn addition, the Higgs invisible branching ratio is then given by\n\\begin{equation}\n\\text{BR}_\\text{inv} = \\frac{\\Gamma_{h \\to ss}}{\\Gamma_{h \\to \\text{SM}} + \\Gamma_{h \\to ss}},\n\\end{equation}\nwhich is constrained to be below $0.24$ at $95\\%$ confidence level \\cite{Khachatryan:2016whc,ATLAS-CONF-2018-031} by direct measurements and below about $0.17$ by statistical fits of Higgs couplings \\cite{Giardino:2013bma,Belanger:2013xza}, which excludes singlet masses below $53$ GeV.\n\nThe portal coupling $\\lambda_{hs}$ is determined by demanding that the correct relic abundance $\\Omega_{c} h^{2} = 0.120$ \\cite{Aghanim:2018eyx} is produced in the freeze-out of the singlet. The PMPC in Eq.~(\\ref{the constraint for our vacuum VS singlet vacuum}) allows us to determine the singlet self-coupling $\\lambda_s$ as a function of the dark matter mass, as is presented in the Fig. \\ref{dark matter self-coupling as a function of dark matter mass}. The self-coupling $\\lambda_s$ becomes non-perturbative as the dark matter mass reaches electroweak scale. \n\nDirect detection also limits \nthe dark matter mass so that the only allowed mass region is near the Higgs resonance. The PMPC thus predicts that the dark matter mass must be in the resonance region. \n\n\\begin{figure*}[tb]\n\t\\begin{center}\n\t\\includegraphics[width=0.45\\linewidth]{RealSingletLowMassPlot.pdf}\n\t\\includegraphics[width=0.45\\linewidth]{RealSingletHighMassPlot.pdf}\n\t\\end{center}\n\t\\caption{\\label{dark matter self-coupling as a function of dark matter mass} The value of $\\mathbb{Z}_2$ real singlet DM self-coupling $\\lambda_s$ as a function of DM mass according to PMPC in Eq. \\eqref{the constraint for our vacuum VS singlet vacuum}. The portal coupling is fixed so that the correct relic abundance is obtained. Left panel: the self-coupling $\\lambda_s$ near the Higgs resonance, $m_s\\in [53,70]$ GeV. Right panel: the self-coupling $\\lambda_s$ away from the resonance, $m_s\\in [65,1000]$ GeV. }\n\\end{figure*}\n\n\n\n\n\n\\section{Model II: $\\mathbb{Z}_2$ pseudo-Goldstone DM} \n\\label{sec:model2}\n\nThis model is an extension of SM with a complex scalar $S$ carrying a global $U(1)$ charge. The most general scalar potential invariant under global $U(1)$ transformation is given by\n\\begin{equation}\\label{}\n\\begin{split}\n V_0 &= \\mu_h^2 \\lvert H \\rvert^{2} + \\mu_s^2 \\lvert S \\rvert^{2} +\\lambda_h \\lvert H \\rvert^{4}\n +\\lambda_s \\lvert S \\rvert^{4} \\\\\n &+\\lambda_{hs} \\lvert H \\rvert^{2} \\lvert S \\rvert^{2},\n\\end{split}\n\\end{equation}\nwhere $H$ is the SM Higgs doublet. \nThe global $U(1)$ symmetry is spontaneously broken as the scalar $S$ acquires a non-zero vacuum expectation value. This would result in a massless Goldstone boson in the physical spectrum. In order to avoid this, the global $U(1)$ symmetry is explicitly broken into its discrete subgroup $\\mathbb{Z}_2$ by the following soft term:\n\\begin{equation}\nV_{\\mathbb{Z}_2}=-\\frac{\\mu^{\\prime 2}}{4}S^2+\\text{h.c},\n\\end{equation}\nso the full potential becomes\n\\begin{equation}\\label{Z2 scalar potential}\nV=V_0+V_{\\mathbb{Z}_2}.\n\\end{equation}\n\nThe model was first considered in \\cite{Barger:2008jx,Chiang:2017nmu}, but only in the Higgs resonance region. The suppression of the direct detection cross section of the pseudo-Goldstone DM was noted in \\cite{Gross:2017dan}; there has been large interest in the phenomenology of the model \\cite{Huitu:2018gbc,Azevedo:2018oxv,Azevedo:2018exj,Alanne:2018zjm,Cline:2019okt,Kannike:2019wsn}.\n\nThe parameter $\\mu^{\\prime 2}$ can be made real and positive by absorbing its phase into the field $S$. When $\\mu^{\\prime 2}$ is positive, then the VEV of $S$ is real. The reality of the VEV can be easily seen by writing the singlet in polar coordinates:\n\\begin{equation}\nS=\\frac{1}{\\sqrt{2}}r e^{i\\theta}.\n\\end{equation}\nThe soft term now becomes\n\\begin{equation}\nV_{\\mathbb{Z}_2}=-\\frac{\\mu^{\\prime 2}}{4}r^2\\cos (2\\theta).\n\\end{equation}\nThe soft term is the only part of the potential that depends on the phase of the singlet $S$. Therefore the global minimum of the potential is the minimum of the soft term. Because $\\mu^{\\prime 2}>0$, the minimum is obtained when $\\cos(2\\theta)=1$. It follows that\n\\begin{equation}\n2\\theta =\\pm 2\\pi n\\Rightarrow \\theta=0,\\pi,\n\\end{equation}\nand therefore the VEV is real.\n\n\nWe parametrize the fields in the unitary gauge as:\n\\begin{equation}\nH=\\frac{1}{\\sqrt{2}}\\left(\\begin{array}{c}\n0\\\\\nh+v_h\n\\end{array}\n\\right),\n\\end{equation}\nand\n\\begin{equation}\nS=\\frac{1}{\\sqrt{2}}(v_s+s+i\\chi).\n\\end{equation}\nThe Higgs vacuum expectation value is $v_h=246$ GeV and $v_s$ is real.\n\nThe scalar potential in Eq. (\\ref{Z2 scalar potential}) is invariant under two different $\\mathbb{Z}_2$ symmetries,\n\\begin{equation}\nS\\to -S,\n\\end{equation}\nand the CP-like\n\\begin{equation} \nS\\to S^\\ast.\n\\end{equation} \nAfter the symmetry breaking the latter symmetry becomes equivalent to $\\chi\\to -\\chi$, which stabilizes $\\chi$, making it a dark matter candidate. \n\nWe will now move on to study the vacuum structure of the potential in Eq. (\\ref{Z2 scalar potential}).\n\n \n\\subsection{Different vacua}\nThe potential in \\eqref{Z2 scalar potential} is bounded from below if\n\\begin{equation} \n\\lambda_h > 0,\\quad \\lambda_s > 0, \\quad 2\\sqrt{\\lambda_h \\lambda_s}+\\lambda_{hs} > 0, \n\\end{equation}\n There are three possible vacuum configurations in field space $(h,s,\\chi)$: the mixed vacuum $(v_h, v_s,0)$, the Higgs vacuum $(v'_h, 0, 0)$ and the singlet vacuum $(0, v'_s, 0)$. Let us study if there can be degenerate minima.\n\n\n\\subsubsection{Our vacuum $(v_h,v_s,0)$}\n\nIn order for the pseudo-Goldstone to be the DM candidate, our Universe has to reside in the mixed vacuum. This vacuum $(v_h,v_s,0)$ satisfies the following minimization conditions:\n\\begin{align}\n\\mu_h^2 \n& = -\\lambda_h v_h^2-\\frac{1}{2}\\lambda_{hs} v_s^2\\label{Our vacuum minimization condition h},\\\\\n\\mu_s^2 & = -\\lambda_s v_s^2-\\frac{1}{2}\\lambda_{hs} v_h^2+\\frac{1}{2}\\mu^{\\prime 2},\n\\label{Our vacuum minimization condition s}\n\\end{align}\nwhich can be inverted in favor of the VEVs:\n\\begin{align}\nv_h^2 & = \\frac{-4\\lambda_s \\mu_h^2 + 2\\lambda_{hs}\\mu_s^2 -\\lambda_{hs}\\mu^{\\prime 2}}{4\\lambda_h\\lambda_s-\\lambda_{hs}^2}>0\\label{solved VEVs 1},\\\\\nv_s^2 & = \\frac{-4\\lambda_h \\mu_s^2 + 2\\lambda_{hs}\\mu_h^2 -2\\lambda_{h}\\mu^{\\prime 2}}{4\\lambda_h\\lambda_s-\\lambda_{hs}^2}>0.\\label{solved VEVs 2}\n\\end{align}\n\nThe mass matrix in the $(h,s,\\chi)$-basis is given by\n\\begin{equation}\n\\mathcal{M}^2=\\begin{pmatrix}\n2\\lambda_h v_h^2 & \\lambda_{hs} v_h v_s & 0\\\\\n\\lambda_{hs} v_h v_s & 2\\lambda_s v_s^2 & 0\\\\\n0 & 0 & \\mu^{\\prime 2}\n\\end{pmatrix}.\n\\end{equation}\n\nThe determinant of the $2\\times 2$ CP-even block matrix has to be strictly positive in order for the mixed configuration to be a physical minimum, yielding the condition\n\\begin{equation}\\label{positive CP-even masses}\n4\\lambda_h\\lambda_s-\\lambda_{hs}^2 >0.\n\\end{equation}\nWe have chosen the parameter $\\mu^{\\prime 2}$ to be positive and therefore all the masses will be positive. The mixed vacuum configuration $(v_h, v_s,0)$ is therefore the minimum to which the other configurations are compared.\n\nWhen one combines Eq. \\eqref{positive CP-even masses} with Eqs. \\eqref{solved VEVs 1} and \\eqref{solved VEVs 2}, the following inequalities are obtained:\n\\begin{align}\n-4\\lambda_s\\mu_h^2+2\\lambda_{hs}\\mu_s^2-\\lambda_{hs}\\mu^{\\prime 2} &>0,\\label{VEV positivity condition 1}\\\\\n-4\\lambda_h\\mu_s^2+2\\lambda_{hs}\\mu_h^2+2\\lambda_h\\mu^{\\prime 2} &>0.\\label{VEV positivity condition 2}\n\\end{align}\n\n\n\\subsubsection{The Higgs vacuum $(v'_h,0,0)$}\n\nWhile the Higgs vacuum $(v'_h,0,0)$ gives rise to correct electroweak symmetry breaking, in this case it is not the pseudo-Goldstone boson but the complex singlet $S$ that is the DM candidate, contrary to the present assumptions. The mass matrix in this vacuum is given by\n\\begin{equation}\n\\mathcal{M}^2=\\begin{pmatrix}\n-2\\mu_h^2 & 0 & 0\\\\\n0 & -\\frac{\\lambda_{hs}\\mu_h^2}{2\\lambda_h}+\\mu_s^2-\\frac{\\mu^{\\prime 2}}{2} & 0\\\\\n0 & 0 & -\\frac{\\lambda_{hs}\\mu_h^2}{2\\lambda_h}+\\mu_s^2+\\frac{\\mu^{\\prime 2}}{2}\n\\end{pmatrix}.\n\\end{equation}\nBy using Eq. (\\ref{VEV positivity condition 2}), one sees that the mass in the middle is negative. Therefore this vacuum configuration is not a minimum.\n\n\\subsubsection{The singlet vacuum $(0,v'_s,0)$}\n\nThe singlet vacuum $(0,v'_s,0)$ can in principle be a minimum. Let us compare, however the potential energy of this minimum to that of our mixed minimum $(v_h,v_s,0)$. Their difference is\n\\begin{equation}\n\\begin{split}\n V(v_h,v_s,0)&-V(0,v'_s,0) =-\\frac{1}{4}\\frac{1}{4\\lambda_h\\lambda_s-\\lambda_{hs}^2} \\frac{1}{\\lambda_s}\n \\\\\n &\\times \\left[\\frac{\\lambda_{hs}}{2}(2\\mu_s^2-\\mu^{\\prime 2})-2\\lambda_s\\mu_s^2\\right]^2<0,\n\\end{split}\n\\end{equation}\naccording to Eq.~(\\ref{positive CP-even masses}). Therefore, we have that\n\\begin{equation}\nV(v_h,v_s,0)0$, the VEV of $S$ is negative.\n\n\nWe now parametrize the scalar fields in unitary gauge as\n\\begin{equation}\nH=\\frac{1}{\\sqrt{2}}\\begin{pmatrix}\n0\\\\\nh+v_h\n\\end{pmatrix},\n\\end{equation}\nand\n\\begin{equation}\nS=\\frac{1}{\\sqrt{2}}(v_s+s+i\\chi).\n\\end{equation}\nThe Higgs vacuum expectation value is $v_h=246$ GeV.\nThe pseudo-scalar $\\chi$ can be dark matter \\cite{Kannike:2019mzk}. The $\\mathbb{Z}_3$-symmetric complex singlet model has also been studied in the phase where the $\\mathbb{Z}_3$ remains unbroken \\cite{Belanger:2012zr,Hektor:2019ote}. \n\nThe scalar potential in Eq. (\\ref{Z3 scalar potential}) is invariant under the CP-like $\\mathbb{Z}_2$-symmetry\n\\begin{equation}\nS\\to S^\\ast.\n\\end{equation}\nThis symmetry is becomes \n\\begin{equation} \n\\chi\\to -\\chi,\n\\end{equation}\nwhen $S$ acquires VEV. This symmetry stabilizes $\\chi$, making it a dark matter candidate. \n\nWe next study the vacuum structure of the model.\n \n\\subsection{Different vacua}\n\nThe potential in Eq. (\\ref{Z3 scalar potential}) is bounded from below, if\n\\begin{equation} \n\\lambda_h > 0,\\quad \\lambda_s > 0, \\quad 2\\sqrt{\\lambda_h \\lambda_s}+\\lambda_{hs} > 0. \n\\end{equation}\n \nThere are three possible vacuum configurations in the field space $(h,s,\\chi)$: the mixed vacuum $(v_h,v_s,0)$, the Higgs vacuum $(v'_h,0,0)$ and the singlet vacuum $(0,v'_s,0)$.\n\n\n\n\n\\subsection{Our vacuum $(v_h,v_s,0)$}\n\nIn order for the pseudo-Goldstone to be the DM candidate, our Universe has to reside in the mixed vacuum. \nThe vacuum $(v_h,v_s,0)$ satisfies the following minimization conditions:\n\\begin{eqnarray}\n\\mu_h^2 \n& = & -\\lambda_h v_h^2-\\frac{1}{2}\\lambda_{hs} v_s^2,\n\\label{Our vacuum minimization condition h}\\\\\n\\mu_s^2 & = &-\\lambda_s v_s^2-\\frac{1}{2}\\lambda_{hs} v_h^2-\\frac{3}{2\\sqrt{2}}v_s{\\mu'_s}.\\label{Our vacuum minimization condition s}\n\\end{eqnarray} \nThe minimization conditions are now hard to invert in favor of VEVs as minimization condition of $\\mu_s^2$ contains a linear term in $v_s$, due to the soft cubic self-coupling of $S$. The VEVs are solved to be:\n\\begin{align}\nv_s & = \\frac{-\\frac{3\\mu'_s}{\\sqrt{2}}\\pm 2\\sqrt{\\frac{1}{\\lambda_h}\\left(4\\lambda_h\\lambda_s-\\lambda_{hs}^2\\right)\\left(\\frac{\\lambda_{hs}\\mu_h^2}{2\\lambda_h}-\\mu_s^2\\right)+\\frac{9\\mu_{s}^{\\prime 2}}{8}}}{\\frac{1}{\\lambda_h}\\left(4\\lambda_h\\lambda_s-\\lambda_{hs}^2\\right)},\n\\label{vs our vacuum minimum}\\\\\nv_h & = \\pm\\frac{\\sqrt{-\\lambda_{hs} v_s^2-2\\mu_h^2}}{\\sqrt{2}\\sqrt{\\lambda_h}}.\n\\end{align}\nWe choose the positive sign $v_s$-solution to be the VEV in our minimum. The negative sign $v_s$-solution is not a minimum as one of the scalar masses is necessarily negative in this case.\n\nThe mass matrix in the $(h,s,\\chi)$-basis is given by\n\\begin{equation}\\label{Z3 our minimum masses}\n\\mathcal{M}^2=\\left(\\begin{array}{ccc}\n2\\lambda_h v_h^2 & \\lambda_{hs} v_h v_s & 0\\\\\n\\lambda_{hs} v_h v_s & 2\\lambda_s v_s^2+\\frac{3}{2\\sqrt{2}}v_{s} \\mu'_s & 0\\\\\n0 & 0 & -\\frac{9v_s \\mu'_s}{2\\sqrt{2}}\n\\end{array}\n\\right).\n\\end{equation}\n\nThe $2 \\times 2$ block in the upper-left corner of Eq. \\eqref{Z3 our minimum masses} can be diagonalized by an orthogonal transformation\n\\begin{equation}\nU^T M^2 U = \\textrm{diag}(m_h^2,m_H^2),\n\\end{equation}\nwhere\n\\begin{equation}\nM^2=\\left(\\begin{array}{cc}\n2\\lambda_h v_h^2 & \\lambda_{hs} v_h v_s \\\\\n\\lambda_{hs} v_h v_s & 2\\lambda_s v_s^2+\\frac{3}{2\\sqrt{2}}vs \\mu'_s \n\\end{array}\n\\right),\n\\end{equation}\nand\n\\begin{equation}\nU=\\left(\\begin{array}{cc}\n\\cos\\theta & \\sin\\theta \\\\\n-\\sin\\theta & \\cos\\theta \n\\end{array}\n\\right).\n\\end{equation}\nThe mixing angle $\\theta$ is defined as\n\\begin{equation}\n\\tan 2\\theta =-\\frac{\\lambda_{hs}v_h v_s}{\\lambda_h v_h^2-\\lambda_s v_s^2-\\frac{3}{4\\sqrt{2}}v_s \\mu'_s}.\n\\end{equation}\nWith the knowledge that we must choose the positive sign solution of $v_s$, the parameters in the potential can be interchanged to more suitable ones. The potential parameters $\\mu_h^2$, $\\mu_s^2$, $\\lambda_h$, $\\lambda_s$, $\\lambda_{hs}$ and $\\mu'_s$ are replaced by $v_h$, $v_s$, $m_h$, $m_H$, $m_\\chi$ and the mixing angle $\\theta$:\n\\begin{align}\n\\mu_h^2 & = -\\frac{1}{4}(m_h^2+m_H^2)+\\frac{1}{4v_h}(m_H^2-m_h^2)\n \\notag\n \\\\\n &\\times (v_h\\cos 2\\theta -v_s\\sin 2\\theta),\\\\\n\\mu_s^2 & = -\\frac{1}{4}(m_h^2+m_H^2)+\\frac{1}{4v_s}(m_h^2-m_H^2)\n\\notag\n\\\\\n&\\times(v_s\\cos 2\\theta +v_h\\sin 2\\theta)+\\frac{1}{6}m_\\chi^2,\\\\\n\\lambda_h & = \\frac{m_h^2+m_H^2+(m_h^2-m_H^2)\\cos 2\\theta}{4 v_h^2},\\\\\n\\lambda_s & = \\frac{3(m_h^2+m_H^2)+3(m_H^2-m_h^2)\\cos 2\\theta+2m_\\chi^2}{12 v_s^2},\\\\\n\\lambda_{hs} & = -\\frac{(m_h^2-m_H^2)\\sin 2\\theta}{2 v_h v_s},\\\\\n\\mu'_s & = -\\frac{2\\sqrt{2}}{9}\\frac{m_\\chi^2}{v_s}.\n\\end{align}\n\n\\subsection{The Higgs vacuum $(v'_h,0,0)$}\n\nWhile the Higgs vacuum $(v'_h,0,0)$ gives rise to correct electroweak symmetry breaking, in this case it is not the pseudo-Goldstone boson but the complex singlet $S$ that is the DM candidate, contrary to the present assumptions.\nThe mass matrix in this vacuum is given by\n\\begin{equation}\n\\mathcal{M}^2=\\begin{pmatrix}\n-2\\mu_h^2 & 0 & 0\\\\\n0 & -\\frac{\\lambda_{hs}\\mu_h^2}{2\\lambda_h}+\\mu_s^2 & 0\\\\\n0 & 0 & -\\frac{\\lambda_{hs}\\mu_h^2}{2\\lambda_h}+\\mu_s^2\n\\end{pmatrix}.\n\\end{equation}\nAll the masses can be positive simultaneously. This minimum can be degenerate with our mixed vacuum, leading into the condition\n\\begin{equation}\n\\mu^{\\prime 2}=\\frac{1}{\\lambda_h}\\left(4\\lambda_h\\lambda_s-\\lambda_{hs}^2\\right)\\left(-\\frac{\\lambda_{hs}}{2\\lambda_h}\\mu_h^2+\\mu^2_s\\right),\n\\end{equation}\nor equivalently,\n\\begin{equation}\\label{higgs only vacuum degeneracy}\nm_\\chi^2 = \\frac{9m_h^2 m_H^2}{m_h^2\\cos^2\\theta+m_H^2\\sin^2\\theta}.\n\\end{equation}\n\n\n\\subsection{The singlet vacuum $(0,v'_s,0)$}\nIn the singlet vacuum $(0,v'_s,0)$, only the complex singlet $S$ gets a VEV and electroweak symmetry is not broken.\nThe minimization condition is given by\n\\begin{equation}\n\\mu_s^2=-\\lambda_s {v'_s}^2-\\frac{3}{2\\sqrt{2}}v'_s \\mu'_s.\n\\end{equation}\nThe VEV $v'_s$ can be found to be\n\\begin{equation}\nv'_s=\\frac{-3\\mu'_s\\pm \\sqrt{-32\\lambda_s\\mu_s^2+9\\mu^{\\prime 2}}}{4\\sqrt{2}\\lambda_s}.\n\\end{equation}\n\nThe mass matrix in this vacuum is given by\n\\begin{equation}\n\\mathcal{M}^2=\\left(\\begin{array}{ccc}\n\\mu_h^2-\\frac{\\lambda_{hs}}{2 \\lambda_s}\\left[\\mu_s^2+\\frac{3\\mu'_s v'_s}{2\\sqrt{2}}\\right] & 0 & 0\\\\\n0 & -2\\mu_s^2-\\frac{3\\mu'_s v'_s}{2\\sqrt{2}} & 0\\\\\n0 & 0 & -\\frac{9\\mu'_s v'_s}{2\\sqrt{2}}\n\\end{array}\n\\right).\n\\end{equation}\n\nThis vacuum configuration does not seem to be able to be a minimum when the mixed vacuum $(v_h,v_s,0)$ and the Higgs vacuum $(v'_h,0,0)$ are degenerate. It is hard to show analytically, but all numerical evidence points to that.\n\n\\begin{figure}[tb]\n\\begin{center}\n\\includegraphics[width=0.5\\textwidth]{direct_detection}\n\\end{center}\n\\caption{Direct detection cross section for $\\mathbb{Z}_{3}$ pGDM. The XENON1T result \\cite{Aprile:2018dbl} is shown in solid line and prediction \\cite{Ni2017} for the XENONnT in dashed line. The red points are excluded by the Higgs invisible width or a too large Higgs total width.}\n\\label{fig:Z3:direct:detection}\n\\end{figure}\n\n\n\n\\subsection{Final verdict on $\\mathbb{Z}_3$ pGDM from the PMPC}\nThe two vacuum configurations $(v_h,v_s,0)$ and $(v'_h,0,0)$ can be minima and degenerate at the same time. This degeneracy of these two minima will impose the condition in Eq. (\\ref{higgs only vacuum degeneracy}), which aproximately reads:\n\\begin{equation}\nm_\\chi\\approx 3 m_H,\n\\end{equation}\nso strictly speaking, $\\chi$ is on the borderline of not being a true pseudo-Goldstone boson.\nThe third vacuum configuration $(0,v'_s,0)$ does not seem to be a minimum when the $(v_h,v_s,0)$ and $(v'_h,0,0)$ are degenerate. This is also checked only numerically. \n\n\nWe also impose constraints from collider experiments. The fit of the Higgs signal strengths constrains the value of the mixing angle to $\\cos^2\\theta\\ge 0.9$~\\cite{Beacham:2019nyx}. If the mass of $x \\equiv \\chi \\text{ or }h_2$ is less than $m_h\/2$, then the Higgs decay width into this particle is given by \\cite{Kannike:2019mzk}\n\\begin{equation}\n \\Gamma_{h \\to x x} = \\frac{g_{hxx}}{8 \\pi} \\sqrt{1 - 4 \\frac{m_x^2}{m_h^2}}\n\\end{equation}\nwith\n\\begin{align}\n g_{h1 \\chi \\chi} &= \\frac{m_h^2 + m_\\chi^2}{v_S},\n \\\\\n g_{h1 h_2 h_2} &= \\frac{1}{v v_S} \\left[ \\left(\\frac{1}{2} m_h^2 + m_2^2\\right) (v \\cos \\theta + v_S \\sin \\theta) \\right.\n \\notag\n \\\\\n & \\left. \n + \\frac{1}{6} v m_\\chi^2 \\cos \\theta \\right] \\sin 2 \\theta.\n\\end{align}\nThe latest measurements of the width of an SM-like Higgs boson\ngives $\\Gamma^{h_1}_\\text{tot}\\,=3.2^{+2.8}_{-2.2}$ MeV, with 95\\% CL limit on\n$\\Gamma_\\text{tot} \\leq 9.16$ MeV \\cite{Sirunyan:2019twz}.\nIn addition, the Higgs invisible branching ratio is then given by\n\\begin{equation}\n\\text{BR}_\\text{inv} = \\frac{\\Gamma_{h \\to \\chi\\chi}}{\\Gamma_{h_1 \\to \\text{SM}} + \\Gamma_{h \\to \\chi\\chi} + \\Gamma_{h \\to h_2 h_2}},\n\\end{equation}\nwhich is constrained to be below $0.24$ at $95\\%$ confidence level \\cite{Khachatryan:2016whc,ATLAS-CONF-2018-031} by direct measurements and below about $0.17$ by statistical fits of Higgs couplings \\cite{Giardino:2013bma,Belanger:2013xza}. Decays of the resulting $h_{2}$ into the SM may be visible \\cite{Huitu:2018gbc}.\n\nA scan of the parameter space with the condition \\eqref{higgs only vacuum degeneracy} imposed is shown in Fig.~\\ref{fig:Z3:direct:detection}. We take the mass of the other CP-even scalar $m_{2} \\in [10, 1500]$ and the mixing angle $\\sin \\theta \\in \\pm [0.01, 0.3]$. The VEV $v_{s}$ is fitted from the DM relic density $\\Omega_{c} h^{2} = 0.120$ \\cite{Aghanim:2018eyx}. The points in red are excluded by the Higgs invisible width and\/or by an overly large Higgs total width.\n\n\\section{Discussion and Conclusion}\n\\label{sec:conc}\n\n\nIn this work we have argued that the principle of multiple point criticality must be generalized to all vacuum states, implying that all minima of a scalar potential \nmust be degenerate. Indeed, if the PMPC follows from some fundamental underlying principle, its validity must not depend on the particulars of the model.\nFor example, in the most studied case, one minimum is taken to be the SM EW Higgs minimum and another one the Planck-scale minimum generated by the running of Higgs quartic coupling. We argue that the principle must apply to all minima in multi-scalar models, whether they are generated radiatively or at tree level.\n\nPhenomenological implications of this formulation of the PMPC are straightforward but profound -- the PMPC becomes a tool to constrain the parameter space of \nmulti-scalar models. Since the scalar singlet models represent one of the most popular class of dark matter models, we exemplified the usefulness of the PMPC by studying how different realizations of scalar DM models are constrained by imposing the PMPC requirements. \n\nIn the simplest possible model, studied in Section~\\ref{sec:model1}, the DM is a real scalar singlet stabilized by a $\\mathbb{Z}_2$ symmetry. The PMPC predicts that the singlet self-coupling hits a Landau pole at DM mass of about $300$~GeV, invalidating the model above that scale. Due to the stringent direct detection constraints, the DM mass in this model is predicted to be near the Higgs resonance.\n\nThe second model under investigation, the $\\mathbb{Z}_2$ symmetric pseudo-Goldstone DM model, studied in Section~\\ref{sec:model2}, cannot have degenerate minima at all. Therefore it cannot be consistent with the PMPC requirements. \n\nThe third model under investigation, the $\\mathbb{Z}_3$ symmetric pseudo-Goldstone DM model, studied in Section~\\ref{sec:model3}, can have degenerate minima. Applying the PMPC in this case leads to a simple relation between the DM and the second scalar masses and excludes DM masses below about $190$~GeV via the measured Higgs width. This model represents an example how the PMPC can constrain the phenomenology of multi-scalar models.\n\nIn conclusion, if the PMPC is indeed realized in Nature, which we have assumed in this work, it may lead to significant constraints and simplifications of the phenomenology of multi-scalar models. Further studies of degenerate tree-level minima in those models is worth of pursuing. \n\n\n\\vspace{5mm}\n\\noindent \\textbf{Acknowledgement.} We thank Gia Dvali, Luca Marzola, \nAlessandro Strumia, Daniele Teresi and Hardi Veerm{\\\"a}e for discussions on the origin of PMPC. \nThis work was supported by the Estonian Research Council grants PRG434, PRG803, PRG803, MOBTT5 and MOBTT86, and by the EU through the European Regional Development Fund CoE program TK133 ``The Dark Side of the Universe.\" \n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{\\bf 1. Introduction.}\n\nIn two previous works, [\\putref{DandK2}], [\\putref{DandK1}], we have\ninvestigated the thermal quantities for scalar and spinor fields on the\nspace--time T$\\times$S$^d$, where $d$ is odd. The conformal scalar\ninternal energy, $E$, for example, is a linear combination of `partial'\nenergies, $\\epsilon_t,\\, t=1,2,\\ldots$, (see below). These quantities possess\na known behaviour under modular transformations and are central to\nelliptic function theory. In particular, a temperature inversion symmetry\ncan be exhibited.\n\nFrom basic elliptic properties, it was shown that $E$ could be expressed\nas a polynomial in $\\epsilon_2$ and $\\epsilon_3$, while, for the specific heat, and\nall higher derivatives, one has to include $\\epsilon_1$, preferably via a\nmodular covariant derivative. The free energy involves an integration and\nthings are not so simple as obstructions arise to simple modular\nbehaviour. It is this aspect that is explored in the present work. I use\nthis thermal angle just to motivate the introduction of various\nquantities, the analysis of which allows us to make contact with various\ntopics in the theory of modular forms. Some of the procedures can be\ngiven a physical terminology, which might be suggestive.\n\n\\section{\\bf 2. Mellin transforms and the period polynomial.}\n\nIt is advantageous to begin from the basic Mellin transform, used {\\it e.g. } by\nMalurkar and Hardy, for the Eisenstein series (or partial internal\nenergy, [\\putref{DandK1}] [\\putref{DandK2}]),\n $$\\eqalign{\n \\epsilon_t({\\mbox{{\\ssb\\char98}}})&=-{B_{2t}\\over 4t}+\n \\sum_{n=1}^\\infty {n^{2t-1}q^{2n}\\over1-q^{2n}}\\cr\n &={1\\over2}\\,\\zeta_R(1-2t)+{1\\over2\\pi i}\\int_{c-i\\infty}^{c+i\\infty}ds\\,\\Gamma(s)\\,\n \\zeta_R(s)\\,\\zeta_R(s-2t+1)\\,(2\\pi{\\mbox{{\\ssb\\char98}}})^{-s}\\,,}\n \\eql{mell1}$$\nwhere $q=e^{i\\pi\\tau}$ and $c$ lies above $2t$. I am using\n${\\mbox{{\\ssb\\char98}}}=-i\\tau= \\beta\/2\\pi=1\/\\xi$ as yet another convenient variable.\nAlthough, physically, $\\mbox{{\\ssb\\char98}}$ is real, if we wish to consider general\nmodular transformations then we must allow it to become complex.\n\nThe Mellin transform has been used systematically in finite temperature\nfield theory and elsewhere in discussions of asymptotic limits. The\nsurvey by Elizalde {\\it et al}, [\\putref{Elizalde2}], is useful and I\nfurther mention, as being somewhat relevant here, Cardy [\\putref{Cardy}],\nKutasov and Larsen, [\\putref{KandL}]. Terras, [\\putref{Terras}] pp.55,229,\ncan be consulted for aspects of the mathematical side. \\marginnote{OTHERS} In\nnumber theory see Hardy and Littlewood, [\\putref{HandL}], and Landau,\n[\\putref{Landau,Landau2}].\n\nEquation (\\puteqn{mell1}) can be looked upon as a continuous expansion in\n${\\mbox{{\\ssb\\char98}}}$. Various limits in ${\\mbox{{\\ssb\\char98}}}$ are uncovered by displacing the\nintegration contour. The pole in the integrand at $s=2t$ corresponds to\nthe Planck term, while the remaining one at $s=0$ corresponds to minus\nthe Casimir value, which is the first term in (\\puteqn{mell1}). The high\ntemperature limit (${\\mbox{{\\ssb\\char98}}}\\to0$) follows by pushing the contour to the\nleft, beyond the Casimir pole, the contribution from which cancels the\nfirst term in (\\puteqn{mell1}) leaving just the effect of the Planck pole\nand a correction that tends to zero exponentially. Conversely, on moving\nthe contour all the way to the right only the zero temperature Casimir\nterm remains.\n\nThis cosmetic asymmetry can be avoided (if desired) by also extracting\nthe Planck term, defining,\n $$\n \\epsilon_t^{\\rm sub}({\\mbox{{\\ssb\\char98}}})=\\epsilon_t(\\mbox{{\\ssb\\char98}})-{1\\over2}\\,\\zeta_R(1-2t)\n \\big(1+({i\\mbox{{\\ssb\\char98}}})^{-2t}\\big)\\,,\n \\eql{rege}$$\nand writing\n $$\\eqalign{\n\\epsilon_t^{\\rm sub}({\\mbox{{\\ssb\\char98}}})={1\\over2\\pi i}\n\\int_{c-i\\infty}^{c+i\\infty}ds\\,\\Gamma(s)\\,\n \\zeta_R(s)\\,\\zeta_R(s-2t+1)\\,(2\\pi{\\mbox{{\\ssb\\char98}}})^{-s}\\,,}\n \\eql{mell2}$$\nwhere now $c$, at the moment, lies anywhere {\\it between} $0$ and $2t$.\n\nThe quantity in the integrand is the Hecke $L$--series of the normalised\nEisenstein modular form,\n $$\n L(G_t,s)=\\zeta_R(s)\\,\\zeta_R(s-2t+1)\\,,\n \\eql{eisl1}$$\nsatisfying the (typical) functional relation\n $$\n (2\\pi)^{-s}\\Gamma(s)\\,L(G_t,s)=(-1)^t(2\\pi)^{s-2t}\\Gamma(2t-s)\\,L(G_t,2t-s)\\,.\n \\eql{gfr}$$\n\nI now consider a multiple integration with respect to ${\\mbox{{\\ssb\\char98}}}$ of the\ninternal energy. Actually for the free energy only one integration is\nneeded but the general case is mathematically important. So integrate\n$\\epsilon_t^{\\rm sub}$, $h$ times to give, using (\\puteqn{mell2}),\n $$\\eqalign{\n{1\\over\\Gamma(h)}&\\int_\\mbox{{\\ssb\\char120}}^\\infty d{\\mbox{{\\ssb\\char98}}}\\,({\\mbox{{\\ssb\\char98}}}-\\mbox{{\\ssb\\char120}})^{h-1}\n\\,\\epsilon_t^{\\rm sub}({\\mbox{{\\ssb\\char98}}})\\cr\n & ={1\\over2\\pi i}\n\\int_{c-i\\infty}^{c+i\\infty}ds\\,{\\Gamma(s-h)\\over(2\\pi)^s}\\,\n \\zeta_R(s)\\,\\zeta_R(s-2t+1)\\,\\mbox{{\\ssb\\char120}}^{-s+h}\\,,}\n \\eql{mint}$$\nwith $h2t-2$ and, continued in $s$, satisfies the\nfunctional relation,\n $$\n (2\\pi)^{-s}\\,\\Gamma(s)\\,L({\\cal F},s)=(-1)^t(2\\pi)^{s-2t}\\,\\Gamma(2t-s)\\,L({\\cal F},2t-s)\\,,\n \\eql{lfr}$$\nby virtue of the (inversion) modularity of ${\\cal F}$, and conversely. This\nis usually attributed to Hecke, [\\putref{Hecke}] (but see the Appendix).\nTerras, [\\putref{Terras}] p.229, gives a useful pedagogical treatment of\nthis topic.\n\nThen the periods could be {\\it defined}, in general, by,\n[\\putref{KandZ}],[\\putref{Zagier}],\n $$\\eqalign{\n r_n({\\cal F})&=\\lim_{s\\to n+1}\\,{\\Gamma(s)\\over(2\\pi)^{s}}\\,L({\\cal F},s)\n \\cr\n &\\equiv\\lim_{s\\to n+1}L^*({\\cal F},s)\\,,\n }\n \\eql{periods2}$$\nbecause this coincides with the integral expression, (\\puteqn{periods}), for\ncusp forms.\n\nThe explicit values of $r_n$ ($n=0,\\ldots 2t-2$) for Eisenstein series\n(which is all that is required because of a decomposition theorem for\nmodular form space) are easily computed from (\\puteqn{eisl1}) and are given\nin [\\putref{KandZ}]. They agree with (\\puteqn{even}) and (\\puteqn{moments}), as\nalready noted.\n\nHowever, the Dirichlet series, $L^*({\\cal F},s)$ has poles at $s=0$ and\n$s=2t$. A means of incorporating these values is given by Zagier,\n[\\putref{Zagier}], using the period polynomial. One notes that the\nsummation in the cusp period polynomials can be extended,\n $$\\eqalign{\n r_{\\cal F}(\\mbox{{\\ssb\\char120}})&=\\sum_{n\\in\\,\\mbox{\\open\\char90}}(-1)^n\\comb{2t-2}{n}\\,r_n({\\cal F})\\,\n \\mbox{{\\ssb\\char120}}^{2t-2-n} \\cr\n&=\\sum_{n\\in\\,\\mbox{\\open\\char90}}(-1)^n\\comb{2t-2}{n}\\,i^{n+1}\\,L^*({\\cal F},n+1)\\,\n\\mbox{{\\ssb\\char120}}^{2t-2-n}\\,,}\n \\eql{ppol6}$$\nbecause the binomial coefficient vanishes outside the range\n$n=0,\\ldots,2t-2$. For {\\it non}--cusp forms, however, the sum extends\nfrom $-1$ to $n=2t-1$, the nonzero end values arising from the\nsingularities of $L^*({\\cal F},s)$ mentioned above. The resulting expression\ncan then be taken as a {\\it definition} of the period polynomial of a\nnon--cusp form.\n\nAgain, for the Eisenstein series, $G_{2t},\\,=\\epsilon_{t}$, the computation is\nstraightforward, [\\putref{Zagier}], and (\\puteqn{ppol6}) with (\\puteqn{eisl1})\nyields, ($n=2j-1$),\n $$\\eqalign{\n r_G(x)={(2t-2)!\\over2(2\\pi)^{2t}}\\bigg(2\\pi\n \\zeta_R&(2t-1)\\,\\big((-1)^{t-1}\\mbox{{\\ssb\\char120}}^{2t-2}-1\\big)-\\cr\n &(2\\pi)^t\\sum_{j=0}^t\\,(-1)^{t-j}{B_{2j}\\,B_{2t-2j}\\over(2j)!(2t-2j)!}\n \\,\\mbox{{\\ssb\\char120}}^{2j-1}\\bigg)\\,,\n }\n \\eql{eperiodp}\n $$\nwhich is $1\/\\mbox{{\\ssb\\char120}}$ times a polynomial. For comparison with\n[\\putref{Zagier}], Zagier's $X$ equals my $i\\mbox{{\\ssb\\char120}}$.\n\nI now comment on this construction. The periods, and period polynomials,\nfor non--cusp forms are defined by analogy to those of cusp forms, the\njustification being that the expressions, for ${\\cal F}(i\\infty)$ equal to\nzero, reduce to those for cusp forms. The means to attain this is not\nunique. Instead of the `zeta--function' regularisation, as in the\ndefinition of $L$, (\\puteqn{ell}), I employ the {\\it subtracted} form, as in\n(\\puteqn{rege}),\n $$\n {\\cal F}^{\\rm sub}\\equiv {\\cal F}-a_0-{a_0\\over\\tau^{2t}}\\,,\\quad\n a_0={\\cal F}(i\\infty)\\,,\n $$\nwhose moments are non--infinite and provide an alternative definition of\nthe periods of a non--cusp form which equally reduces to that for cusp\nforms when $a_0=0$. This subtraction maintains inversion modularity, but\nviolates translational, whereas the subtraction of just $a_0$ does the\nreverse. Making use of my earlier considerations of Eisenstein series,\nsee (\\puteqn{ppol}), I define the subtracted period polynomial as an\nintegral,\n $$\n r_{\\cal F}^{\\rm sub}(\\mbox{{\\ssb\\char120}})=\\int_0^{\\infty}d\\mbox{{\\ssb\\char98}}\\,\n (\\mbox{{\\ssb\\char98}}-\\mbox{{\\ssb\\char120}})^{2t-2}{\\cal F}^{\\rm sub}(\\mbox{{\\ssb\\char98}})\\,.\n $$\n\nFor Eisenstein series, (\\puteqn{ppol}) shows that $r_G^{\\rm sub}(\\mbox{{\\ssb\\char120}})$ is\nproportional to $\\overline P_t(\\mbox{{\\ssb\\char120}})$ whose explicit form is given in\n(\\puteqn{ppoly2}).\n\nIt is necessary to compare this definition with that of Zagier,\n(\\puteqn{eperiodp}). The difference is the absence of the `end point' terms\nproportional to $1\/\\mbox{{\\ssb\\char120}}$ and $\\mbox{{\\ssb\\char120}}^{2t-1}$. This can be remedied by\nincluding, in the terminology of this paper, the regularised\ncontributions from the subtracted terms, the Casimir and Planck terms,\nwhich can be accomplished simply by expanding the closed contour in\n(\\puteqn{ppoly2}) to include the corresponding poles at $s=0$ and $s=2t$.\nThis just extends the sum by the two end points giving the well known\nquantity,\n $$\\eqalign{\n\\overline R_t(\\mbox{{\\ssb\\char120}})=2\\pi\\,\\zeta_R(2t-1)\\big((i \\mbox{{\\ssb\\char120}})^{2t-2}-1\\big)-4i\n \\sum_{j=0}^t\\zeta_R(2j)\\,\\zeta_R(2t-2j)(i \\mbox{{\\ssb\\char120}})^{2t-2j-1}\n \\,,}\n \\eql{ppoly3}\n $$\nwhich agrees, up to a factor, with Zagier's expression given in\n(\\puteqn{eperiodp}).\n\nFurthermore, the use of the subtracted form, ${\\cal F}^{\\rm sub}$, allows one\nto give a cleaner integral expression for the enlarged period polynomial\nthan that in [\\putref{Zagier}], {\\it i.e. },\n $$\n r_{\\cal F}^{\\rm enl}(\\mbox{{\\ssb\\char120}})=\\int_0^{\\infty}d\\mbox{{\\ssb\\char98}}\\,\n (\\mbox{{\\ssb\\char98}}-\\mbox{{\\ssb\\char120}})^{2t-2}{\\cal F}^{\\rm sub}(\\mbox{{\\ssb\\char98}})+\n {a_0\\over2t-1}\\big(\\mbox{{\\ssb\\char120}}^{2t-1}+\\mbox{{\\ssb\\char120}}^{-1}\\big)\\,.\n $$\n\nWhile equivalent to the method in [\\putref{Zagier}], I believe the contour\napproach is neater and is capable of being extended to the general theory\nof periods. Zagier, [\\putref{Zagier}], shows that logistic simplifications\noccur in this theory when it is extended to include non--cusp forms such\nas the Eisenstein series.\n\nAs an historical point, it will be noticed that the above expressions\nhave been available since the time of Ramanujan, with contour derivations\nby Malurkar, [\\putref{Malurkar}], and Guinand, [\\putref{Guinand}]. In order\nto investigate this point further, I return to the mode sum expressions\nfor the thermodynamic related quantities.\n\n\\section{ \\bf 3. Back to summations.}\nI start with the Eisenstein series forms and work the development to\nparallel the quantities of the previous section. This allows me to\nintroduce the useful calculations of Smart, [\\putref{Smart}], on the\nEpstein function (which will arise later) using its relation with\nEisenstein series.\n\nThe subtracted partial energy $\\epsilon_t^{\\rm sub}$ in (\\puteqn{rege})\ncorresponds to the the `regularised' Eisenstein series,\n $$\\eqalign{\n G_t^{\\,\\rm sub}(\\tau)&=\\sumdasht{m=-\\infty}{\\infty}\\,\\sumdasht{n=-\\infty}{\\infty}\n {1\\over (m\\tau+n)^{2t}}\\cr\n &=\\,G_t(\\tau)-2\\zeta_R(2t)(1+\\tau^{-2t})\\,,}\n \\eql{geesub}$$\nThe relation is, [\\putref{DandK1}],\n $$\n \\epsilon^{\\rm sub}_t(\\tau)={(2t-1)!\\over2(2\\pi i)^{2t}}\\, \\overline\n G^{\\,\\rm sub}_t(\\tau)\\,.\n \\eql{inten5}$$\n\nNow, corresponding to (\\puteqn{mint2}), integrate (\\puteqn{geesub}) with\nrespect to $\\tau$ from $\\tau$ to $i\\infty$, $2t-1$ times to give,\nformally,\n $$\n {1\\over(2t-1)!}\\,\\phi_{2t}(\\tau)\\,,\n \\eql{geeint}$$\nwhere $\\phi$ is defined by, [\\putref{Smart}], eqn.(2.11),\n $$\n \\phi_{2t}(\\tau)=\\sumdasht{m=-\\infty}{\\infty}\\,\\sumdasht{n=-\\infty}{\\infty}\n {1\\over m^{2t-1}}{1\\over m\\tau+n}\\,.\n \\eql{phis}$$\n\nThe behaviour of the series $\\phi_{2k}(\\tau)$ under modular\ntransformations, in particular $\\tau\\to-1\/\\tau$, has been especially\nconsidered by Smart, [\\putref{Smart}], who refers to the $\\phi$'s as `{\\it\nmodular forms of weight $2-2t$ with rational period functions}'.\n\nUsing the partial fraction identity (see also Hurwitz, [\\putref{Hurwitz}],\nEisenstein, [\\putref{Eisenstein}]),\n $$\n {1\\over m^{2t-1}}{1\\over m\\tau+n}=\n -{\\tau^{2t-1}\\over n^{2t-1}}{1\\over m\\tau+n}+\\sum_{j=1}^{2t-1}(-1)^{j+1}\n {\\tau^{2t-j-1}\\over m^j\\,n^{2t-j}}\\,,\n \\eql{pf}$$\none finds easily the inversion relation, [\\putref{Smart}], eqn.(1.10b),\n $$\n \\phi_{2t}(\\tau)-\\tau^{2t-2}\\phi_{2t}(-1\/\\tau)=-4\\sum_{j=1}^{t-1}\n \\tau^{2t-2j-1}\\zeta_R(2j)\\,\\zeta_R(2t-2j)\\,.\n \\eql{sm1}$$\nAlso, under translations,\n $$\n \\phi_{2t}(\\tau+1)-\\phi_{2t}(\\tau)=2{\\zeta_R(2t)\\over\\tau(\\tau+1)}\\,.\n \\eql{sm2}$$\n\nUp to a simple factor, one might expect $\\phi_{2t}$ to be the same as\n${\\overline\\phi}_{2t}$ in (\\puteqn{mint2}). However the cocycle functions in\n(\\puteqn{sm1}) and (\\puteqn{ppoly2}) differ by the two {\\it even} powers. I\nnote that these are related to the $j=1$ and $j=2t-1$ terms in the sum in\n(\\puteqn{pf}) which have gone out in the passage to the difference,\n(\\puteqn{sm1}), by antisymmetry. However the summation over $m$,\nrespectively $n$, for these values of $j$ is conditionally convergent and\nis therefore subject to ambiguity, that is to say, a choice has to be\nmade, as in the discussion of the Eisenstein series, $G_2$. The\nregularised Mellin transform approach results in a different definition\nof the series (\\puteqn{phis}) to that used by Smart, [\\putref{Smart}],p.3. The\nrelation is easily found by taking the even powers in (\\puteqn{ppoly2}) over\nto the left and then we see that one has the relation, (further remarks\nare given later),\n $$\n {\\overline\\phi}_{2t}(\\mbox{{\\ssb\\char120}})=\ni\\phi_{2t}(\\tau)-2\\pi \\zeta_R(2t-1)\\,,\\quad \\mbox{{\\ssb\\char120}}=-i\\tau\\,,\n \\eql{connec}\n $$\nalso agreeing with (\\puteqn{trans}) and (\\puteqn{sm2}).\n\nSmart, [\\putref{Smart}], also gives the relation with a known Lambert\nseries as\n $$\\eqalign{\n S_t(\\tau)\\equiv\\sum_{m=1}^\\infty{1\\over m^{2t-1}}{q^{2m}\\over1-q^{2m}}\n &=-{1\\over4\\pi i}\\,\\phi_{2t}(\\tau)-{1\\over2\\pi i\\tau}\\,\n\\zeta_R(2t)-{1\\over2} \\zeta_R(2t-1)\\cr\n&={1\\over4\\pi}\\,\\big(\\overline\\phi_{2t}(\\mbox{{\\ssb\\char120}})+{2\\over\\mbox{{\\ssb\\char120}}}\\zeta_R(2t)\\big)\\cr\n&\\equiv{1\\over4\\pi}\\,\\overline\\psi_{2t}(\\mbox{{\\ssb\\char120}})\\,,\n}\n \\eql{ls}$$\nwhere $\\overline\\psi_{2t}$ is the function having $\\overline R_{2t}(\\mbox{{\\ssb\\char120}})$,\n(\\puteqn{ppoly3}), as its inversion cocycle function,\n $$\n {\\overline\\psi}_{2t}(\\mbox{{\\ssb\\char120}})-(i\\mbox{{\\ssb\\char120}})^{2t-2}\\,{\\overline\\psi}_{2t}(1\/\\mbox{{\\ssb\\char120}})\n= \\overline R_{2t}(\\mbox{{\\ssb\\char120}})\\,,\n \\eql{are1}$$\nas is easily confirmed.\n\nThe introduction of $\\overline\\psi$ has restored translational\ninvariance, (the Planck term has been put back in (\\puteqn{ls})),\n $$\n \\overline\\psi_{2t}(\\mbox{{\\ssb\\char120}}-i)- \\overline\\psi_{2t}(\\mbox{{\\ssb\\char120}})=0\\,.\n \\eql{per}$$\nIt is almost convention to consider modular integrals, for example, to be\nperiodic ({\\it e.g. } Knopp, [\\putref{Knopp}]). See some later comments.\n\nEquations (\\puteqn{are1}) and (\\puteqn{ls}) imply an inversion relation for\nthe series $S_t$ which has a certain history, some of which is detailed\nby Berndt, [\\putref{Berndt}], p.153, who says that it was first written\ndown by Ramanujan. As is clear, it is quite the same as the expression\nfor the Eisenstein period polynomials.\n\nThe Lambert series, $S_t(\\tau)$, has been considered by Apostol,\n[\\putref{Apostol}], Berndt, [\\putref{Berndt}], Grosswald, [\\putref{Grosswald}],\nand Smart [\\putref{Smart}], for example. It can be related to statistical\nmechanics in the following way.\n\nOne can resum the statistical expression for the partial free energy in a\nstandard manner, beginning, ($q=e^{\\pi i\\tau}=e^{-\\pi\/\\xi}$),\n $$\n f_t=\\epsilon_{t,0}-{\\xi\\over2\\pi}\\sum_{n=1}^\\infty n^{2t-2}\\log(1-q^{2n})\n =\\epsilon_{t,0}-{\\xi\\over2\\pi}\\sum_{m,n=1}^\\infty{n^{2t-2}\\over m}q^{2mn}\\,.\n \\eql{ssum}$$\n\nThe scaled entropy is defined by,\n $$\n s_t(\\xi)={\\partial f_t(\\xi)\\over\\partial \\xi}=-{1\\over2\\pi}\\bigg(1-{2\\pi\\over\\xi}D\n\\bigg) \\sum_{m,n=1}^\\infty{n^{2t-2}\\over m}q^{2mn}\\,.\n \\eql{scent}$$\n\nConcentrate now on the summation and rewrite it,\n $$\n \\sum_{m,n=1}^\\infty{n^{2t-2}\\over m}q^{2mn}=D^{2t-2}\n\\sum_{m=1}^\\infty{1\\over m^{2t-1}}{q^{2m}\\over1-q^{2m}}\\,,\n \\eql{sum9}$$\nin terms of $S_t(\\tau)$, (\\puteqn{ls}).\n\nThis Lambert series has been discussed in connection with generalised\nDedekind $\\eta$--functions by Berndt, [\\putref{Berndt}], and it is worth\nnoting that, in a certain sense, $\\log \\eta$ is a modular form of zero\nweight, in agreement with the results $\\epsilon_1=-D\\log\\eta$,\n[\\putref{DandK2}], which could therefore be written,\n $\n \\epsilon_1=-{\\cal D}\\log\\eta\\,,\n $\nin terms of the covariant derivative, ${\\cal D}$.\n\nThe question, not answered here, is whether there is an `expansion' or\nrepresentation theorem for modular {\\it integrals} combining that for\nmodular forms (of {\\it any} real weight, see [\\putref{Rankin}]). It is not\nlikely that one can find a pure elliptic formula. As evidence I cite the\nparticular evaluations of Smart who gives at the lemniscate point\n($\\tau=i,\\mbox{{\\ssb\\char120}}=1$),\n $$\n \\overline\\psi_{4}(1)={7\\pi^4\\over90}-2\\pi\\zeta_R(3)\\,,\n $$\nand, more generally, if $2t=0({\\rm mod} 4)$,\n $$\n S_t(i)={\\zeta_R(2t)\\over2\\pi}-{\\zeta_R(2t-1)\\over2}+{1\\over2\\pi}\n \\sum_{j=1}^{t-1}(-1)^{j+1}\\zeta_R(2t-2j)\\zeta_R(2j)\\,,\n $$\na result actually due to Lerch which follows from (\\puteqn{are1}) with\n(\\puteqn{ppoly3}).\n\nAs a final general comment in this section, it is always possible to\nobtain the required quantities as (infinite) $q$-series by direct\nintegration, but these do not count as closed forms. For example, one\nintegration yields the $q$-series for the partial free energy in terms of\nthe arithmetic functions, $\\sigma_k(m)$, (the sum of the $k$-th powers of\nthe divisors of $m$),\n $$\n f_t=\\epsilon_{t,0}+{1\\over\\beta}\\sum_{m=1}^\\infty\\sigma_{2t-1}(m)\\,{q^{2m}\\over m}\\,,\n $$\nwhich amounts only to doing the $n$-summation in ({\\puteqn{ssum}), or an\nintegration of the standard relation, (\\puteqn{mell1}),\n $$\n \\epsilon_t=-{B_{2t}\\over4t}+\\sum_{m=1}^\\infty \\sigma_{2t-1}(m)\\,q^{2m}\\,.\n $$\n\nOne also has the formula, {\\it e.g. } [\\putref{Rad,Glaisher}],\n $$\n S_t(\\tau)=\\sum_{m=1}^\\infty \\sigma_{-2t+1}(m)\\,q^{2m}=\n \\sum_{m=1}^\\infty {\\sigma_{2t-1}(m)\\over m^{2t-1}}\\,q^{2m}\\,,\\quad\\forall\\, t\\,,\n \\eql{Ssi}$$\nin agreement with the result of a $(2t-1)$-fold integration.\n\\section{\\bf 4. Another approach via modular integrals.}\nThe period functions are cocycle functions associated with Eichler\ncohomology and modular integrals. These notions put a different slant on\nthe preceding formulae. In effect one begins again.\n\nThe best place to start is the important (but particular) analytical result\nof Bol, [\\putref{Bol}], which states that, for {\\it any} reasonable function,\n$\\varphi(\\tau)$,\n $$\n \\big(D^{r+1}\\varphi\\big)(\\gamma\\tau)=\n (c\\tau+d)^{r+2}D^{r+1}\\big((c\\tau+d)^r\\,\\varphi(\\gamma\\tau)\\big)\\,,\n \\eql{bol}$$\nwhere $r$ is a nonnegative number, and I use,\n $$(D^{r+1}\\varphi)(\\tau)\\equiv \\bigg(q{d\\over dq}\\bigg)^{r+1}\\varphi(\\tau)\\,,\n $$\nwith,\n $$\n \\gamma=\\left(\\matrix{*&*\\cr\n c&d}\\right)\\in {\\rm SL}(2,\\mbox{\\open\\char82})\\,.\n $$\nThis shows that, {\\it if} $\\varphi$ is an automorphic form of weight\n$-r$, then $D^{r+1}\\,\\varphi$ is an automorphic form of weight $2+r$, see\nalso Petersson, [\\putref{Petersson}]. Note that it is the ordinary\nderivative that enters here. The proof of (\\puteqn{bol}) follows by\ninduction.\n\nThe classic discussion of {\\it modular integrals} can proceed by asking\nfor the inverse of Bol's result. That is, given a form of weight $2+r$,\n$\\rho(\\tau)$, as a source, can one integrate the differential equation,\n $$\n D^{r+1}\\varphi=\\rho\n \\eql{de}$$\nto obtain $\\varphi$, a form of weight $-r$? Clearly this will not be so\nin general because of the freedom introduced by constants of integration.\nThat is, one would expect $\\varphi$ to behave as a form of weight $-r$ up\nto an additional piece, the general form of which can be determined from\nthe differential equation (\\puteqn{de}) and Bol's theorem, (\\puteqn{bol}).\n\nGenerally we would write (ignoring any multiplier systems for\nsimplicity),\n $$\n \\varphi(\\gamma\\tau)=(c\\tau+d)^{-r}\\big(\\varphi(\\tau)+P(\\gamma,\\tau)\\big)\\,,\n \\eql{co2}$$\nwhere $P(\\gamma,\\tau)$ has to be a polynomial in $\\tau$ of order $\\le r$,\nThe inverse to Bol's result is that, if any $\\varphi$ satisfies\n(\\puteqn{co2}), then $\\rho$, of (\\puteqn{de}), is a form of weight $(2+r)$.\nThis follows most simply by substituting (\\puteqn{co2}) into (\\puteqn{bol}).\n\n$P$ must satisfy the 1--cocycle consistency condition, [\\putref{Eichler}],\n $$\\eqalign{\n P(\\gamma_1\\gamma_2,\\tau)&=(c_2\\tau+d_2)^r\\,P(\\gamma_1,\\gamma_2\\tau)+P(\\gamma_\n 2,\\tau)\\cr\n &\\equiv P(\\gamma_1,\\tau)\\big|\\gamma_2+P(\\gamma_2,\\tau)\\,.\n }\n \\eql{cocycle}$$\n\n$\\varphi$ is known as an {\\it automorphic} or {\\it Eichler integral}, of\nweight $-r$, and $P$ is the {\\it associated period polynomial},\n[\\putref{Eichler,HandK}]. I have also introduced the useful, standard\n`stroke' operator, $\\big|\\gamma$, sometimes denoted $\\big|[\\gamma]$.\n\nThe {\\it general} solution (equivalently the indefinite integral) of the\ndifferential equation (\\puteqn{de}) is a particular integral plus the\ncomplementary function (a zero mode) {\\it i.e. }\n $$\n \\varphi(\\tau)={(2\\pi i)^{r+1}\\over\\Gamma(r+1)}\n \\int_{\\tau_0}^\\tau d\\sigma\\, (\\tau-\\sigma)^{r}\\,\\rho(\\sigma)+\\Theta(\\tau)\\,,\n \\eql{gensol}\n $$\nwhere $\\Theta(\\tau)$ is a polynomial of degree $r$. The integral, in the\nupper half plane, is path independent if $\\rho$ has weakish analytic\nproperties and, if desired, the solution can be made independent of the\nlower limit by suitable adjustment of $\\Theta(\\tau)$. In elementary fashion,\n$\\Theta(\\tau)$ is determined by the first $r$ derivatives of $\\varphi(\\tau)$\nat $\\tau_0$.\n\nAll this is general analysis. In modular applications to SL(2,$\\mbox{\\open\\char90}$),\n$\\tau_0$ is chosen as the cusp $i\\infty$ and the limits in (\\puteqn{gensol})\nreversed. In terms of Fourier expansions, a holomorphic cusp form source\nis\n $$\n \\rho(\\tau)=\\sum_{m=1}^\\infty \\rho_m\\,q^{2m}\\,,\n \\eql{fe1}$$\nand the particular Eichler integral of $\\rho$, is, from (\\puteqn{gensol}),\n $$\n \\varphi(\\tau)=\\sum_{m=1}^\\infty {\\rho_m\\over m^{r+1}}\\,q^{2m}\\,,\n \\eql{fe2}$$\nwhich has the same periodicity as $\\rho$, a result of choosing\n$\\tau_0=i\\infty$.\n\nThe problem in the present Eisenstein case has already been noted. The\nEisenstein series, $G_t(\\tau)$ does not vanish at $i\\infty$ (or $0$). It\nis a `non--cusp' form and the direct integration of (\\puteqn{de}) meets\nobstacles at an elementary level. This can be appreciated simply from the\nappearance of $1\/m$ factors during integration, and $m$ can be zero in\n$G_t$. This fits in because the $m=0$ terms in $G_t$ are those making\n$G_t$ nonzero at $i\\infty$, the zero temperature value. Saying it yet\nagain, the Fourier transform of $G_t$ begins with an $m=0$ term, the\nCasimir contribution, so (\\puteqn{fe2}) makes no immediate sense.\n\nOne way of avoiding these obstructions, is to regularise $G_t$ by\ndropping the offending term(s) as in (\\puteqn{geesub}) and consider the\ndifferential equation,\n $$\\eqalign{\n D^{2t-1}\\,\\varphi_{2t}=-4\\pi\\epsilon_t^{\\rm sub}\\,,}\n \\eql{gee}$$\nwhich has the particular integral, $\\varphi_{2t}=\\overline\\phi_{2t}$\nwhere,\n $$\n \\overline\\phi_{2t}(\\tau)=-4\\pi{(2\\pi i)^{2t-1}\\over\\Gamma(2t-1)}\n \\int_\\tau^{i\\infty}d\\tau'\\,(\\tau'-\\tau)^{2t-2}\n \\,\\epsilon_t^{\\rm sub}(\\tau')\\,,\n $$\nthe constants and boundary conditions being chosen so as to reproduce\n(\\puteqn{mint2}). The upper limit is a cusp. Smart's function, $\\phi_{2t}$,\nresults if one adds a complementary constant (and multiplies by $i$).\n(See (\\puteqn{connec}).)\n\nBecause of the subtraction, the quantity $\\epsilon_t^{\\rm sub}(\\tau)$ is not\nexactly a modular form, rather one has, expressed slightly generally,\n($r=2t-2$),\n $$\n \\epsilon(\\gamma\\tau)=(c\\tau+d)^{r+2}\\big(\\epsilon(\\tau)+E(\\gamma,\\tau)\\big)\n \\eql{submod}\n $$\nand also\n $$\n \\phi(\\gamma\\tau)=(c\\tau+d)^{-r}\\big(\\phi(\\tau)+P(\\gamma,\\tau)\\big)\\,,\n \\eql{co3}\n $$\nwhich are related by the same differential equation, (\\puteqn{de}),\n $$\n D^{r+1}\\phi=-4\\pi\\epsilon\\,.\n $$\nBol's formula, (\\puteqn{bol}), relates $P$ and $E$,\n $$\n D^{r+1}P=-4\\pi E\\,,\n \\eql{bol2}\n $$\nand these relations generalise the usual ones applicable to cusp forms.\nThe cocycle formula (\\puteqn{cocycle}) still holds of course, as (\\puteqn{co3})\nis the same as (\\puteqn{co2}).\n\nI remark that $\\overline\\phi$ and $\\epsilon^{\\rm sub}$, not being periodic, do\nnot quite possess Fourier expansions, as in (\\puteqn{fe1}) and (\\puteqn{fe2}).\nHowever $\\overline\\psi$, in view of (\\puteqn{per}), does,\n $$\n \\overline\\psi(\\mbox{{\\ssb\\char120}})=\\sum_{m=1}^\\infty \\psi_m\\,q^{2m}\\,,\n $$\n\nBecause $P$, this time, satisfies (\\puteqn{bol2}), it is not required to be\na polynomial. Generally, since the modular group is generated \\footnote{\nBeware that some authors, {\\it e.g. } Knopp, swap the usage of the symbols $S$\nand $T$.} by $S$ and $T$, $P(\\gamma,\\tau)$ is a {\\it rational period\nfunction}. In the present case, however, $P(S,\\tau)$ {\\it is} a\npolynomial, given by (\\puteqn{ppoly2}),\n $$\n P(S,\\tau)=\\overline P_t(\\tau)\\,,\n \\eql{pt}$$\nwhile the (non--zero) $P(T,\\tau)$, which follows from (\\puteqn{trans}),\n $$\n P(T,\\tau)=2\\,\\zeta_R(2t)\\bigg({1\\over \\tau}-{1\\over\\tau+1}\\bigg)\\,,\n \\eql{ps}$$\nis a ratio of polynomials. Of course, $P(\\pm{\\bf 1},\\tau)=0$ and I remark\nthat the polynomial $P(S,\\tau)$ runs from $\\tau^0$ to $\\tau^{2t-2}$.\n\nAs a simple check of the algebra, set $\\gamma_1=\\gamma_2=S$ in the cocycle\nrelation, (\\puteqn{cocycle}). This gives,\n $$\n P(S,\\tau)\\big|(1+S)\\equiv P(S,\\tau)+\\tau^{r}P(S,S\\tau)=\n P(S,\\tau)+\\tau^{r}P(S,-1\/\\tau)=0\n $$\nwhich, for $r=2t-2$, agrees with the explicit form, (\\puteqn{ppoly2}).\n\nAlso, as examples,\n $$\\eqalign{\n P(TS,\\tau)=&2\\,\\zeta_R(2t){\\tau^{2t}\\over1-\\tau}+\\overline P_t(\\tau)\\cr\n P(ST,\\tau)=&{2\\,\\zeta_R(2t)\\over\\tau(1+\\tau)}+\\overline P_t(1+\\tau)\\,.}\n $$\nI note that the period function appears to have poles at $-1$, $0$, $+1$\nand $i\\infty$.\n\nFrom (\\puteqn{submod}) and the form of $\\epsilon^{\\rm sub}$, (\\puteqn{rege}), one\nfinds\n $$\n E(S,\\tau)=0\\,,\\quad E(T,\\tau)=\\bigg({1\\over(\\tau+1)^{2t}}\n -{1\\over\\tau^{2t}}\\bigg)\n $$\nwhich are consistent with (\\puteqn{ps}) and (\\puteqn{pt}).\n\nFor completeness, and to aid comparison, this is a convenient place to\noutline some standard elementary things about period polynomials, mostly\nfor SL$(2,\\mbox{\\open\\char90})$. Two basic `polynomials' are the inversion one,\n$P(S,\\tau)$, and the translation one, $P(T,\\tau)$. For cusp forms, the\nlatter vanishes because of periodicity, $P(T,\\tau)=0$, [\\putref{Knopp}],\ncorresponding to the periodicity of the modular integrals, {\\it cf }\n(\\puteqn{fe1}), (\\puteqn{fe2}). In this case, $P(S,\\tau)$ satisfies the\nEichler--Shimura relations\n $$\\eqalign{\n P(S,\\tau)\\big|(1+S)&=0 \\cr\n P(S,\\tau)\\big|(1+TS+(TS)^2)&=0\n }\n \\eql{ES}\n $$\nwith $S^2=1$ and $(TS)^3=1$. A standard way of proving these uses the\ncomplex integral form for $P$, {\\it e.g. } [\\putref{Zagier}], but they can also\nfollow, algebraically from the cocycle relation as I now show. (See also\nKnopp, [\\putref{Knopp3}], for slightly different details.)\n\nFor the first identity, just set, as above, $\\gamma_1=1$ and $\\gamma_2=S$ in\n(\\puteqn{cocycle}) and it falls straight out. The second identity is\nslightly more work. Using (\\puteqn{cocycle}) twice, one has, dropping the\n$\\tau$ argument temporarily and noting $(TS)^2=(TS)^{-1}$,\n $$\\eqalign{\n P(S)\\big|\\big(1+&TS+(TS)^2\\big)=\\cr\n &P(S)+P(STS)-P(TS)+P(SS^{-1}T^{-1})-P(TSTS)\\,.\n }\n \\eql{ES2}\n $$\n\nFurther, setting $\\gamma_1=T$ in the cocycle condition, and using the\nassumption that $P(T)=0$, yields\n $$\n P(T\\gamma)=P(\\gamma)\\,.\n \\eql{pt0}\n $$\nOne then sees that the first and third terms on the right--hand side of\n(\\puteqn{ES2}) cancel, as do the second and fifth. The fourth term vanishes\non its own because (\\puteqn{pt0}) shows that $P(T^{-1})=0$, and (\\puteqn{ES})\nfollows.\n\nIn one development, the notion of modular integral has been somewhat\ndivorced from that of integration. One says that a {\\it modular integral}\nis any function (with appropriate analyticity properties) that obeys the\nquasi--modular relation (\\puteqn{co2}), where the {\\it period function} $P$\nmust still, of course, satisfy the cocycle condition (\\puteqn{cocycle}).\nMotivation is then provided to look for modular integrals with {\\it\nrational} period functions, {\\it i.e. } ratios of polynomials. In this case it\nseems to be conventional to {\\it assume} that the translational period\nfunction, $P(T,\\tau)$, vanishes, [\\putref{Knopp}], so that both\nEichler--Shimura relations, ({\\puteqn{ES}), still hold. There is a certain\ninterest in computing such rational period functions and investigating\ntheir zeros and poles. The easiest introduction to this topic is by\nKnopp, [\\putref{Knopp3}]. It is, however, somewhat peripheral to our\nconsiderations which are more concerned with non--cusp forms {\\it i.e. } the\nEisenstein series, {\\it e.g. }\\ [\\putref{Zagier}].\n\n\\section{\\bf 5. Green function distributional description.}\n\nDeveloping the formalism a little further, it is possible to give the\ndiscussion of Eichler integrals a distributional or Green function look.\nFor this purpose, it is better to use the `real' variable form and consider\nthe differential equation,\n $$\n {d^{2t-1}\\over d{\\mbox{{\\ssb\\char120}}}^{2t-1}}\\,\\psi(\\mbox{{\\ssb\\char120}})=\\rho(\\mbox{{\\ssb\\char120}})\\,,\n \\eql{rde}$$\nfor $\\mbox{{\\ssb\\char120}}>0$ with solution vanishing at $\\infty$, {\\it cf } (\\puteqn{mint2}),\n $$\n \\psi(\\mbox{{\\ssb\\char120}})=\n\\int_{0^-}^\\infty d{\\mbox{{\\ssb\\char120}}'}\\,\\Phi_{2t-1}(\\mbox{{\\ssb\\char120}}'-\\mbox{{\\ssb\\char120}})\\,\\rho({\\mbox{{\\ssb\\char120}}'})\n \\eql{soln}$$\nwhere the generalised function, $\\Phi_\\alpha(\\mbox{{\\ssb\\char120}})$, is\n $$\n\\Phi_\\alpha(\\mbox{{\\ssb\\char120}})={\\mbox{{\\ssb\\char120}}^{\\alpha-1}_+\\over\\Gamma(\\alpha)}\\,.\n \\eql{genp}$$\nThe {\\it generalised} function $\\mbox{{\\ssb\\char120}}^\\alpha_+$ is concentrated on the positive\n$\\mbox{{\\ssb\\char120}}$-axis, {\\it i.e. } $\\mbox{{\\ssb\\char120}}^\\alpha_+$ is equal to $\\mbox{{\\ssb\\char120}}^\\alpha$ for $\\mbox{{\\ssb\\char120}}\\ge0$ and is\nzero for $\\mbox{{\\ssb\\char120}}<0$, {\\it e.g. } Gelfand and Shilov, [\\putref{GandS}] I,\\S5.5.\n$\\Phi_{2t-1}$ acts as a Green function satisfying,\n $$\n {d^{2t-1}\\over d{\\mbox{{\\ssb\\char120}}}^{2t-1}}\\,\\Phi_{2t-1}(\\mbox{{\\ssb\\char120}}'-\\mbox{{\\ssb\\char120}})=\\delta(\\mbox{{\\ssb\\char120}}'-\\mbox{{\\ssb\\char120}})\\,.\n $$\nThe formal, distributional equivalent is obtained from the convolution, valid\nfor all $\\alpha$ and $\\beta$,\n $$\n \\Phi_\\alpha*\\Phi_\\beta=\\Phi_{\\alpha+\\beta}\\,,\n \\eql{conv}$$\nby setting $\\alpha=-\\beta=2t-1$ and noting\n $$\n\\Phi_{-k}(\\mbox{{\\ssb\\char120}})=\\delta^{(k)}(\\mbox{{\\ssb\\char120}})\\,,\\quad k=0,1,2,\\ldots.\n \\eql{deriv}$$\n\nAn important formula is the Fourier transform of $\\Phi$, [\\putref{GandS}]\nI.p.360,\n $$\n \\int_0^\\infty d\\mbox{{\\ssb\\char120}}\\, \\Phi_\\alpha(\\mbox{{\\ssb\\char120}})\\,e^{i\\sigma \\mbox{{\\ssb\\char120}}}\n ={1\\over (\\sigma+i0)^\\alpha}\\,e^{i\\pi\\alpha\/2}\\,,\n \\eql{ftphi}$$\nand the inverse, which is a (continuous) eigenfunction expansion,\n $$\n \\Phi_\\alpha(\\mbox{{\\ssb\\char120}})=e^{i\\pi\\alpha\/2}{1\\over2\\pi}\\int_{-\\infty}^\\infty d\\sigma\\,\n {e^{-i\\sigma \\mbox{{\\ssb\\char120}}}\\over(\\sigma+i0)^\\alpha}\\,.\n \\eql{invft}$$\n\nIn the present case, the source $\\rho$ is assumed to be periodic under\n$\\mbox{{\\ssb\\char120}}\\to\\mbox{{\\ssb\\char120}}-i$ (the Planck term is not removed) and to converge as\n$\\mbox{{\\ssb\\char120}}\\to\\infty$. So one has the `Fourier' series,\n $$\n \\rho(\\mbox{{\\ssb\\char120}})=\\sum_{m=1}^\\infty \\rho_m\\, e^{-2m\\pi\\mbox{{\\ssb\\char120}}}\\,,\\quad \\mbox{{\\ssb\\char120}}>0\\,.\n \\eql{ftrho}$$\nSubstitution of this and (\\puteqn{invft}) into (\\puteqn{soln}) gives,\n $$\n \\psi(\\mbox{{\\ssb\\char120}})=(-1)^{t+1}{1\\over2\\pi}\\sum_{m=1}^\\infty\\rho_m\n \\int_{-\\infty}^\\infty d\\sigma\\,\n {e^{i\\sigma\\mbox{{\\ssb\\char120}}}\\over(\\sigma+i0)^{2t-1}(\\sigma-2m\\pi i)}\\,.\n \\eql{psi1}\n $$\nClosing the contour in the upper half $\\sigma$--plane I obtain, $\\mbox{{\\ssb\\char120}}>0$,\n $$\n \\psi(\\mbox{{\\ssb\\char120}})={1\\over(2\\pi)^{2t-1}}\\sum_{m=1}^\\infty\n {\\rho_m\\over m^{2t-1}}\\,e^{-2m\\pi\\mbox{{\\ssb\\char120}}}\\,,\n \\eql{soln2}$$\nwhich, it is no surprise to see, is just (\\puteqn{fe2}) derived by a much\nlonger route but one on which we encounter some handy information. It is\nseen also that the source periodicity has been transmitted to the solution.\n\nIncluding an $m=0$ term in the Fourier series, (\\puteqn{ftrho}), still gives\nconvergence at infinity, but of course the solution, (\\puteqn{soln2}),\nbreaks down, as has already been mentioned. As with all Green function\napproaches, when this happens the corresponding eigenfunction has to be\nremoved from the source. This can be done formally by redefining the\nGreen function to exclude this function which is, in the present case, a\nconstant zero mode. Altering the contour in the Fourier integral\n(\\puteqn{invft}) gives a modified Green function,\n $$\n \\Phi^{\\rm sub}_{2t-1}(\\mbox{{\\ssb\\char120}})\\equiv(-1)^{t}{1\\over2\\pi i}\n \\int_{-\\infty+i\\delta}^{\\infty+i\\delta} d\\sigma\\,\n {e^{-i\\sigma\\mbox{{\\ssb\\char120}}}\\over\\sigma^{2t-1}}\\,,\\quad 0<\\delta<2\\,,\n \\eql{phisub}$$\nwhich automatically takes out any constant part of the source. This can\nbe seen by going back to (\\puteqn{psi1}) and replacing the integral by one\nfrom $-\\infty+i\\delta$ to $\\infty+i\\delta$,\n $$\n \\psi(\\mbox{{\\ssb\\char120}})=(-1)^{t+1}{1\\over2\\pi}\\sum_{m=1}^\\infty\\rho_m\n \\int_{-\\infty+i\\delta}^{\\infty+i\\delta} d\\sigma\\,\n {e^{i\\sigma\\mbox{{\\ssb\\char120}}}\\over\\sigma^{2t-1}(\\sigma-2m\\pi i)}\\,.\n \\eql{psi2}\n $$\nFor $ 0<\\delta<2$, if $\\mbox{{\\ssb\\char120}}>0$, closing the contour in the upper half plane\nstill yields (\\puteqn{soln2}). However now, even if the source Fourier\nseries, (\\puteqn{ftrho}), is extended to the constant term, $m=0$, the\nresult is again (\\puteqn{soln2}). This should be denoted by $\\psi^{\\rm\nsub}(\\mbox{{\\ssb\\char120}})$ and one can write\n $$\n \\psi^{\\rm sub}(\\mbox{{\\ssb\\char120}})=\n\\int_{0^-}^\\infty d{\\mbox{{\\ssb\\char120}}'}\\,\\Phi^{\\rm\nsub}_{2t-1}(\\mbox{{\\ssb\\char120}}'-\\mbox{{\\ssb\\char120}})\\,\\rho({\\mbox{{\\ssb\\char120}}'})\\,.\n \\eql{soln3}$$\n\nRetaining the general structure, (\\puteqn{invft}), for any $\\alpha$, one can\ndefine for any $\\alpha$,\n $$\n \\Phi^{\\rm sub}_\\alpha(\\mbox{{\\ssb\\char120}})\\equiv e^{i\\pi\\alpha\/2}{1\\over2\\pi}\n\\int_{-\\infty+i\\delta}^{\\infty+i\\delta}d\\sigma\\,\n {e^{-i\\sigma \\mbox{{\\ssb\\char120}}}\\over\\sigma^\\alpha}\\,.\n \\eql{invft2}$$\nThen $\\Phi^{\\rm sub}_\\alpha(\\mbox{{\\ssb\\char120}})$ is still concentrated on the positive\n$\\mbox{{\\ssb\\char120}}$--axis and the convolution (\\puteqn{conv}) also remains valid,\n $$\n \\Phi^{\\rm sub}_\\alpha* \\Phi^{\\rm sub}_\\beta =\\Phi^{\\rm sub}_{\\alpha+\\beta}\n \\eql{conv2}\n $$\ntogether with (\\puteqn{deriv}),\n $$\n \\Phi^{\\rm sub}_{-k}=\\delta^{(k)}\\,.\n \\eql{deriv2}\n $$\n\nTo check these statements, firstly, if $\\mbox{{\\ssb\\char120}}<0$, closing the contour in\n(\\puteqn{invft2}) in the upper half plane yields zero as the only\nsingularity is at $\\sigma=0$. Then\n $$\\eqalign{\n \\big(\\Phi^{\\rm sub}_\\alpha* \\Phi^{\\rm sub}_\\beta\\big) (\\mbox{{\\ssb\\char120}})=\n {e^{i\\pi(\\alpha+\\beta)\/2}\\over(2\\pi)^2}\\int_{-\\infty}^\\infty d\\mbox{{\\ssb\\char120}}'\n\\int_{-\\infty+i\\delta}^{\\infty+i\\delta}d\\sigma\\,\\int_{-\\infty+i\\delta}^{\\infty+i\\delta}d\\sigma'\\,\n {e^{-i\\sigma\\mbox{{\\ssb\\char120}}-i\\mbox{{\\ssb\\char120}}'(\\sigma'-\\sigma)}\\over\\sigma^\\alpha\\sigma'^\\beta}\n }\n $$\nwhere the lower limit on the $\\mbox{{\\ssb\\char120}}'$ integral is allowed because\n$\\Phi_\\beta^{\\rm sub}(\\mbox{{\\ssb\\char120}}')$ is concentrated on the positive $\\mbox{{\\ssb\\char120}}'$ axis.\n\nSetting $\\sigma=s+i\\delta$, $\\sigma'=s'+i\\delta$, I get, in detail,\n $$\\eqalign{\n \\big(\\Phi^{\\rm sub}_\\alpha* \\Phi^{\\rm sub}_\\beta\\big) (\\mbox{{\\ssb\\char120}})&= e^{i\\pi(\\alpha+\\beta)\/2}{1\\over(2\\pi)^2}\\int_{-\\infty}^\\infty d\\mbox{{\\ssb\\char120}}'\n\\int_{-\\infty}^{\\infty}ds\\,\\int_{-\\infty}^{\\infty}ds'\\,\n {e^{-i\\sigma\\mbox{{\\ssb\\char120}}-i\\mbox{{\\ssb\\char120}}'(s'-s)}\\over\\sigma^\\alpha\\sigma'^\\beta}\\cr\n &= e^{i\\pi(\\alpha+\\beta)\/2}{1\\over2\\pi}\n\\int_{-\\infty}^{\\infty}ds\\,\\int_{-\\infty}^{\\infty}ds'\\,\n {e^{-i\\sigma\\mbox{{\\ssb\\char120}}}\\over\\sigma^\\alpha\\sigma'^\\beta}\\,\\delta(s-s')\\cr\n &=e^{i\\pi(\\alpha+\\beta)\/2}{1\\over2\\pi}\n\\int_{-\\infty}^{\\infty}ds\\,\n {e^{-i\\sigma\\mbox{{\\ssb\\char120}}}\\over\\sigma^\\alpha\\sigma^\\beta}=e^{i\\pi(\\alpha+\\beta)\/2}{1\\over2\\pi}\n\\int_{-\\infty+i\\delta}^{\\infty+i\\delta}d\\sigma\\,\n {e^{-i\\sigma\\mbox{{\\ssb\\char120}}}\\over\\sigma^{\\alpha+\\beta}}\\cr\n &=\\Phi^{\\rm sub}_{\\alpha+\\beta}(\\mbox{{\\ssb\\char120}})\\,.\n }\n $$\n\n Equation (\\puteqn{deriv2}) follows either from (\\puteqn{conv2}) or\n (\\puteqn{invft2}) since the integrand has no singularity, and so one can\n set $\\delta=0$.\n\n\\subsection{\\it Eichler cohomology.}\nAs a side remark, the general solution to (\\puteqn{rde}) consists of the\nparticular integral, (\\puteqn{soln}), plus a solution of the homogeneous\nequation, {\\it i.e. } a zero mode, which in this case is a polynomial,\n$\\Theta(\\mbox{{\\ssb\\char120}})$, of degree at most $2t-2$. The contribution of this to the\ncocycle polynomial, $P$ of (\\puteqn{co2}), constitutes an element of the set\nof {\\it coboundaries} and a cohomology of polynomials can be set up\n(Eichler cohomology). Working with the cohomology classes removes the\nconstants of integration ambiguity in the solution of (\\puteqn{rde}). That\nis, there is a one-to-one correspondence between $\\rho$ and an element of\nthe cohomology group denoted $H^1(\\Gamma,\\Pi_r)$, where $\\Pi_r$ is the\nvector space of polynomials of degree $\\le r=2t-2$ (for cusp forms) and\n$\\Gamma$ is (here) the modular group. The development of these ideas, while\nsignificant generally, is not required here.\n\n\n\\section{\\bf 6. Epstein approach to thermal quantities.}\n\nAn `earlier' form of the free energy, or effective action, follows from\nthe thermal $\\zeta$--function, which, in this case is related to the Epstein $\\zeta$--function.\nStarting from this I will be able to make contact with the Eisenstein\nformulation and derive useful relations. Using the Epstein function is\nequivalent to the non-holomorphic Eisenstein series. As a rule the\nanalysis is harder. In reality one is doing twice the necessary work.\n\nIn some ways we are going backwards. The use of holomorphic forms, {\\it e.g. }\n(\\puteqn{inten5}), should be considered a simplification available through\nworking with an explicit square root of the Laplacian.\n\nAgain splitting up the degeneracy, one is led to define a `partial'\n$\\zeta$--function,\n $$\n \\zeta_t(s,\\beta)={i\\over2\\beta}\\sumdasht{{m,n=\\atop-\\infty}}\n{\\infty}{n^{2t-2}\\over\n \\big(4\\pi^2m^2\/\\beta^2+n^2\/a^2\\big)^s}\\,,\n \\eql{epthzf}$$\nand the related free energy,\n $$\n \\overline F_t=-{\\xi\\over8\\pi}\\lim_{s\\to0}{1\\over s}\\,\n \\sumdasht{{m,n=\\atop-\\infty}}{\\infty}{n^{2t-2}\\over\n \\big(m^2+n^2\\xi^{-2}\\big)^s}\\,.\n \\eql{eff}$$\n\nI have rescaled the free energy by defining, as above and as in\n[\\putref{DandK1}], $\\overline F_t=a F_t$ and have also taken a preliminary\nlimit of $s\\to0$ in an extracted factor of $(2\\pi\/\\beta)^{2s}$. This will\nbe justified in a moment.\n $$\n \\overline F_t=(-1)^{t-1}{\\xi\\over8\\pi}\\bigg({d\\over d(1\/\\xi^2)}\\bigg)^{t-1}\n \\lim_{s\\to0}{\\Gamma(s-t+1)\\over\\Gamma( s+1)}\\,\n \\sumdasht{{m,n=\\atop-\\infty}}{\\infty}{1\\over\n \\big(m^2+n^2\\xi^{-2}\\big)^{s-t+1}}\\,.\n \\eql{eff2}$$\n\nNot surprisingly, the degeneracy factor is replaced by a derivative and the\nresulting sum is precisely an Epstein $\\zeta$--function. This has been used before {\\it e.g. }\n[\\putref{Kennedy}]. The point is that there is no pole at $s=0$, which\njustifies the limit. In terms of the Epstein function, $Z_2$,\n $$\n \\overline F_t=(-1)^{t-1}{\\xi\\over8\\pi}\\bigg({d\\over d(1\/\\xi^2)}\\bigg)^{t-1}\n \\lim_{s\\to0}\\,{\\Gamma(s-t+1)\\over\\Gamma( s+1)}\\,Z_2(2s-2t+2,A)\\,,\n \\eql{eff3}$$\nwhere $A$ is the matrix (the `modulus'),\n $$\n A=\\left(\\matrix{1&0\\cr\n 0&\\xi^{-2}}\\right)\\,.\n $$\nThe idea now is to use the standard functional relation satisfied by\n$Z_2(s,A)$, repeated here for convenience,\n $$\n Z_2(2s,A)={\\pi^{2s-1}\\over\\sqrt{{\\rm det\\,} A}}{\\Gamma(1-s)\\over\\Gamma(s)}\n \\,Z_2(2-2s,A^{-1})\\,,\n \\eql{freln}$$\nto give\n $$\n \\overline F_t={\\xi\\over8\\pi}\\bigg(-{d\\over d(1\/\\xi^2)}\\bigg)^{t-1}\\xi\n \\lim_{s\\to0}\\pi^{2s-2t+1}\\,{\\Gamma(t-s)\\over\\Gamma( s+1)}\\,Z_2(2t-2s,A^{-1})\\,.\n $$\n$Z_2(2s,A)$ has a pole at $s=1$ only, so working at $t>1$ to begin with,\none has\n $$\n \\overline F_t=\\xi{(-1)^{t-1}\\Gamma(t)\\over8\\pi^{2t}}\n \\bigg({d\\over d(1\/\\xi^2)}\\bigg)^{t-1}\\xi\n \\,\\,Z_2(2t,A^{-1})\\,.\n \\eql{epder}$$\nAs the simplest case, set $t=2$ giving the three-sphere. Then\n $$\n \\overline F_3={\\xi\\over8\\pi^4}{d\\over d(1\/\\xi^2)}\\,\\xi\\,Z_2(4,A^{-1})=\n -{\\xi^4\\over16\\pi^4}{d\\over d\\xi}\\,\\xi\\,Z_2(4,A^{-1})\\,,\n \\eql{3free1}$$\nagreeing with the result in [\\putref{Kennedy}], allowing for the changes in\nnotation.\n\nThis formula can be taken further as shown by Epstein [\\putref{Epstein}]\np.633 who gives,\n $$\n Z_2(4, A^{-1})={\\pi^4\\over45}+\\pi\\xi^{-3}\\zeta_R(3)+2\\pi\\xi^{-3}\\chi(q')\n +4\\pi^2\\xi^{-2}\\,D'\\,\\chi(q')\\,,\n \\eql{epz2}$$\nwhere $q'=e^{-\\pi\\xi}$ and\n $$\n \\chi(q')=S_2(-1\/\\tau)=\\sum_{m=1}^\\infty {q'^{2m}\\over m^3(1-q'^{2m})}\\,.\n $$\n\nEquation (\\puteqn{3free1}) with (\\puteqn{epz2}) is appropriate for the high\ntemperature limit since $q'$ is the `inverse' temperature Boltzmann factor.\nThe low temperature form follows by simple homogeneous scaling of the\nEpstein $\\zeta$--function, which gives,\n $$\n \\,Z_2(2s,A^{-1})=\\xi^{-2s}Z_2(2s,A),\n \\eql{scale}$$\nand therefore the equivalent form,\n $$\n Z_2(4, A^{-1})={\\pi^4\\over45}\\xi^{-4}+\\pi\\xi^{-1}\\zeta_R(3)+2\\pi\\xi^{-1}\\chi(q)\n +4\\pi^2\\xi^{-2}\\,D\\,\\chi(q)\\,,\n \\eql{epz21}$$\nwhere $q=e^{-\\pi\/\\xi}$.\n\nThe low temperature form (\\puteqn{epz21}) yields a somewhat simpler expression\nfor the free energy, (\\puteqn{3free1}), as the $\\zeta_R(3)$ term goes out and\nthere is some cancellation that produces,\n $$\n \\overline F_3={1\\over240}+{\\xi\\over2\\pi}D^2\\,\\chi(q)\\,.\n $$\nI have thus regained, at some length, the summation form (\\puteqn{ssum})\nwith (\\puteqn{sum9}).\n\nAll that has been done is to rederive, in a special case and in a\ndetailed way, the standard statistical mode sum, the general form of\nwhich was obtained, in this fashion, in [\\putref{DandK,Dow2}] and we are no\nfurther forward in finding a closed form for the free energy, which is\none of my aims.\n\nThe equivalence of (\\puteqn{epz2}) and (\\puteqn{epz21}), derived here rather\ntrivially (granted Epstein's calculation), is a known inversion identity,\nand can also be derived from (\\puteqn{sm1}). The exact identity is written\nout in Katayama [\\putref{Kata}]. His derivation seems to be similar to that\nvia the Epstein function. Clearly the higher odd sphere expressions will\ninvolve the Riemann $\\zeta$--function\\ at positive odd integers. Katayama writes out\nthe one appropriate for the (partial) five-sphere.\n\nSmart, [\\putref{Smart}], also derives these identities, which is not\nsurprising since his work is concerned with the evaluation of the Epstein\n$\\zeta$--function\\ at integer arguments, through which he is led to the forms\n$\\phi_{2k}$. We see that Epstein had already arrived at the same\ndevelopment. Consult also Bodendiek and Halbritter, [\\putref{BandH}], for a\nrelated treatment.\n\n\\section{\\bf 7. The Epstein--Kober--Selberg--Chowla formula.}\nSmart makes use of the, oft quoted, paper of Selberg and Chowla,\n[\\putref{SandC}], which is concerned, partly, with the evaluation of the\nEpstein function via the Kronecker limit formula \\footnote{Landau,\n[\\putref{Landau2}], gives some interesting history of this famous formula.}\nand elegantly yields some of the finite forms for the complete elliptic\nintegral used in [\\putref{DandK1}]. The important expression for the\nEpstein function quoted by Selberg and Chowla, [\\putref{SandC2}], and\nderived by them in [\\putref{SandC}], (equ.(6)), in terms of $K$--Bessel\nfunctions, is essentially the same as that already given by Epstein,\n[\\putref{Epstein}], p.631, equn.(12), p.622, equn.(16). Rankin\n[\\putref{Rankin3}] obtained the same result and it also occurs in Bateman\nand Grosswald, [\\putref{BandG}]. These are the standard references. See\nalso Terras, [\\putref{Terras}] p.209. Smart, [\\putref{Smart}], employs this\nformula to split the Epstein function into (simpler) holomorphic and\nanti--holomorphic parts using the explicit form of the Bessel function,\n$K_{t-1\/2}$, which henceforth disappears from his analysis. Other\nexpressions were given by Deuring [\\putref{Deuring}] and Mordell,\n[\\putref{Mordell}].\n\nAn early, and little known, derivation is that by Kober, [\\putref{Kober}].\nFollowing, and generalising, Watson, [\\putref{Watson}], he obtains, {\\it\nstarting} from the Bessel function end, the formula (I retain his\nnotation),\n $$\\eqalign{\n {1\\over\\sqrt u}\\,&\\Delta^{(2w+1)\/4}{\\Gamma(w+1\/2)\\over8\\pi^{w+1\/2}}Z_2(a,b,c;w+1\/2)\\cr\n&={u^{-w}\\over4}{\\Gamma(w)\\over\\pi^{w}}\\,\\zeta_R(2w)\n+{u^{w}\\over4}{\\Gamma(w+1\/2)\\over\\pi^{w+1\/2}}\\,\\zeta_R(2w+1)+\\,\\cr\n &\\hspace{************}\n \\sum_{n=1}^\\infty \\sigma_{2w}(n)n^{-w}\\cos(2\\pi vn)K_{w}(2\\pi un)\\,,}\n \\eql{kob1}$$\nwhere the quadratic form in the Epstein function, $Z_2$, is\n$am_1^2+2bm_1m_2+cm_2^2$ and\n $$\n \\Delta=ac-b^2\\equiv a^2\\,u^2\\,,\\quad v\\equiv{b\\over a}\\,.\n $$\nEquation (\\puteqn{kob1}) is identical to the corresponding one in\n[\\putref{Epstein}] and [\\putref{SandC}] and therefore, with justice, should\nbe called the Epstein--Kober formula. It is the {\\it Fourier expansion}\nof the Epstein function in $v$. For the sum of squares case, $v=0$, the\nresult is rederived by Guinand [\\putref{Guinand2}] in the same way.\n\nKober also discusses behaviour under `reciprocal' transformations,\n$u\\to1\/u$. In particular the diagonal case, $v=0$, when one can write\n$u=\\mbox{{\\ssb\\char120}}$.\n\nIn this way of doing things, the Bessel function can be eliminated using\nthe standard formula\n $$\n \\sum_{n=1}^\\infty {\\sigma_{2w}(n)\\over n^{s+w}}=\\zeta_R(s+w)\\,\\zeta_R(s-w)\\,,\n \\eql{sig}$$\nand the Mellin transform (Heaviside's integral),\n $$\n {1\\over4}\\pi^{-s}\\Gamma\\big({s+w\\over2}\\big)\\,\\Gamma\\big({s-w\\over2}\\big)=\n \\int_0^\\infty dy \\,y^{s-1}\\,K_w(2\\pi y)\\,,\\quad {\\rm Re\\,} s>1+|{\\rm Re\\,} w|\\,,\n \\eql{hi}$$\nso that,\n $$\n \\sum_{n=1}^\\infty {\\sigma_{2w}(n)\\over n^w} K_{w}(2\\pi un)\n=\n{1\\over2\\pi i}\\int_{c-i\\infty}^{c+i\\infty}ds\\,u^{-s}\\xi(s+w)\\,\\xi(s-w)\\,,\n \\eql{epmell}$$\nwith $c>1+w$ and where,\n $$\\xi(2w)\\equiv{1\\over2}\\pi^{-w}\\,\\Gamma(w)\\,\\zeta_R(2w)\\,.\n $$\nThis is the same as Deuring's expression, [\\putref{Deuring}]\np.589.\n\nThe integrand possesses a functional equation which follows most easily\nfrom that for the Riemann $\\zeta$--function, $\\xi(2w)=\\xi(1-2w)$, and is,\n $$\n \\xi(s+w)\\,\\xi(s-w)\\equiv f_w(s)=f_w(1-s)\\,.\n \\eql{ffrel}$$\n\nUsing (\\puteqn{ffrel}) one could {\\it derive} the functional equation for the\nEpstein $\\zeta$--function\\ in (\\puteqn{kob1}). The usual derivation employs the inversion\nrelation for generalised $\\theta$--functions, [\\putref{Epstein,Krazer}].\n\nFurthermore, it is possible to reverse the argument and use (\\puteqn{epmell})\nto obtain (\\puteqn{kob1}) (for $v=0$). The diagonal Epstein function is\nparticularly easy to deal with being the $\\zeta$--function\\ on the torus,\nS$^1\\times$S$^1$. It can be expressed very simply in terms of the $\\zeta$--functions\\ on\nthe circle factors (these are Riemann $\\zeta$--functions\\ ) and the result is exactly\n(\\puteqn{kob1}) with (\\puteqn{epmell}) {\\it without} the intervention of the\nBessel function. The individual Riemann $\\zeta$--functions\\ in (\\puteqn{kob1}) arise from\nadjusting the zero modes. All this is standard and can easily be\ngeneralised to higher dimensions.\n\n\nA natural step is to proceed as with the holomorphic modular forms and ask\nfor the corresponding period polynomials which will follow directly from\n(\\puteqn{epmell}) and (\\puteqn{ffrel}) in the standard manner.\n\nThere are four poles in the relevant strip, $-c\\le s\\le c$, and\na straightforward contour calculation produces the result,\n $$\\eqalign{\n \\sum_{n=1}^\\infty {\\sigma_{2w}(n)\\over n^w}K_{w}(2\\pi n u)-&\n {1\\over u}\\sum_{n=1}^\\infty {\\sigma_{2w}(n)\\over n^w}K_{w}(2\\pi n\/u)\\cr\n &={1\\over2}\\xi(2w)(u^{w-1}-u^{-w})+{1\\over2}\\xi(-2w)(u^{-w-1}-u^{w})\\,,}\n \\eql{ppk}$$\nThe right hand side of (\\puteqn{ppk}) might be termed a period function.\n\nThis formula was obtained by Guinand, [\\putref{Guinand2}], but using the\ncomplete equation (\\puteqn{kob1}), with $v=0$, and properties of the Epstein\nfunction. This is unnecessary as I have just shown.\n\nIn order to make a connection with the period polynomials encountered\nearlier,{\\it e.g. } (\\puteqn{are1}), one introduces the formula for the $K$-Bessel\nfunction, written as,\n $$\n K_{t-1\/2}(x)=\\sqrt{2\\over \\pi}\\,x^{t-1\/2}\\,(-1)^{t-1}\\,\\bigg({1\\over x}\n{d\\over dx}\\bigg)^{t-1}{e^{-x}\\over x}\\,,\\quad t\\in\\mbox{\\open\\char90}\\,.\n \\eql{kbess}$$\nThen, somewhat similarly to [\\putref{Smart}],\n $$\\eqalign{\n \\sum_{n=1}^\\infty {\\sigma_{2t-1}(n)\\over n^{t-1\/2}}K_{t-1\/2}(2\\pi n u)&=\n \\sqrt{2\\over\\pi}{(-1)^{t-1}\\over(2\\pi)^{t-1\/2}}\\bigg({1\\over u}\n {d\\over du}\\bigg)^{t-1}\n \\sum_{n=1}^\\infty{\\sigma_{2t-1}(n)\\over n^{2t-1}}\\,q^{2n}\\cr\n &=\\sqrt{2\\over\\pi}{(-1)^{t-1}\\over(2\\pi)^{t-1\/2}}\n \\bigg({1\\over u}{d\\over du}\\bigg)^{t-1}\n S_t(iu)\\,,}\n $$\nso that, for the left--hand side of (\\puteqn{ppk}), one finds,\n $$\n (-1)^{t-1}{1\\over{\\pi}^t}\\bigg(\\bigg({d\\over du^2}\\bigg)^{t-1}\n S_t(iu)-{1\\over u}\\bigg({d\\over d(1\/u^2)}\\bigg)^{t-1}S_t(i\/u)\\bigg)\\,.\n \\eql{krelp} $$\nIt is not obvious from (\\puteqn{are1}) that this reduces so nicely to\n(\\puteqn{ppk}), and I leave this as a question mark.\n\nFurther expressions involving Bessel functions appear in the Appendix.\n\\begin{ignore}\n\nThis can be turned around in the following way. Write (\\puteqn{kbess})\nexplicitly in the form,\\marginnote{THis is wrong}\n $$\n e^{-x}\\sum_{j=0}^{t-1}\n {(t-1+j)!\\over j!(t-1-j)!}\\,(2x)^{-j-1\/2}=\\pi^{-1\/2}\\,K_{t-1\/2}(x)\\,,\n $$\nand solve this for $(2x)^{-j-1\/2}$, $0\\le j\\le t-1$,\n $$\n e^{-x}\\,(2x)^{-j-1\/2}={1\\over\\sqrt\\pi}\\sum_{t=0}^{j-1}{()!\\over()!}\n \\,K_{t-1\/2}(x)\\,.\n $$\nNow set $x=2\\pi n u$ and get\n $$\n {e^{-2\\pi nu}\\over n^{j+1\/2}}=(2\\pi)^{j+1\/2}\n {1\\over\\sqrt\\pi}\\sum_{t=0}^{j-1}{()!\\over()!}\\,K_{t-1\/2}(2\\pi nu)\n $$\n\\end{ignore}\n\\begin{ignore}\n\\section{\\bf More hyperbolic sums.}\nIn terms of the hyperbolic expression, (\\puteqn{speen}), the quantity\ncorresponding to the $c_t$ is,\n $$\n {\\partial\\over\\partial\\beta}\\,E(\\beta)=-{1\\over2^{d+1}}\\sum_{m=1}^\\infty m\\bigg({(d-1)^2\n\\cosh(m\\beta\/2)\\over\\sinh^{d-1}(m\\beta\/2)}+{d(d+1)\\cosh(m\\beta\/2)\n\\over\\sinh^{d+1}(m\\beta\/2)}\\bigg)\\,,\n \\eql{speen2}$$\nwhich means that it is also possible to evaluate the sums,\n $$\n \\sum_{m=1}^\\infty {m^{2p+1}\\cosh(mx)\\over\\sinh^d(mx)}\n $$\nand\n $$\n \\sum_{m=1}^\\infty {m^{2p}\\over\\sinh^d(mx)}\\,,\n $$\nat special values of $x$. Cases with $d=1$ have been considered by Zucker,\n[\\putref{zucker2}] and Ling, [\\putref{Ling2}] and can follow from expansion of\nthe dn function.\n\\section{\\bf The fermion case. More to be done.}\n\nIn [\\putref{DandK1}] it was shown that just as for bosons, the fermion\ninternal energies on all spheres can be expressed in terms of just two\nquantities which could be taken to be the first two partial energies,\n$\\eta_1$ and $\\eta_2$. This was proved without recourse to the boson result\nand in a different manner. The recursion is not similar to (\\puteqn{recurs})\nbut has a cubic form. One might therefore reasonably expect an equation\ncorresponding to (\\puteqn{shr}) to exist, but in a different form and derived\nby a different route, possibly from the Jacobi function ds $u$.\n\nA preliminary fact, that it is useful to confirm, starts from the\nexpression for the partial spinor energy on the circle given in\n[\\putref{DandK1}],\n $$\n \\eta_1(\\beta)=-{1\\over24} \\mbox{{\\goth\\char75}}^2\\,(1-2k^2)\\,,\n $$\nwhich, using (\\puteqn{diffrel}), can be integrated to give,\n $$\n \\int d\\beta\\, \\eta_1(\\beta)={1\\over12}\\int d\\kappa {1-2\\kappa\\over\\kappa(1-\\kappa)}=\n {1\\over12}\\log(\\kappa\\ka'\/16)\\,,\n $$\nwhere $\\kappa=k^2$ and a constant of integration has been chosen to make the\nfree energy and internal energies agree at absolute zero.\n\nI have, in this roundabout way, derived another of Jacobi's formulae,\n $$\n q^{1\/24}\\prod_{j=0}^\\infty(1+q^{2j+1})=2^{1\/6}(kk')^{-1\/12}\\,,\n $$\nthe logarithm of which gives the mode sum form of the circle fermion\nfree energy, as is well known in various contexts.\n\nWe can also see the difficulty if we attempt the same approach for $\\eta_2$\n $$\\eqalign{\n -\\int d\\beta\\,\\eta_2(\\beta)&={1\\over120}\\int d\\kappa\\,\\mbox{{\\goth\\char75}}^2(\\kappa)\\,\n{(7+8\\kappa-8\\kappa^2)\\over\\kappa(1-\\kappa)}\\cr\n&={1\\over120}\\int d\\kappa\\,\\mbox{{\\goth\\char75}}^2(\\kappa)\\,\\bigg({7\\over\\kappa}\n+{7\\over\\kappa'}-8\\bigg)\\,,}\n $$\nand the integral is not obvious. I merely note the hypergeometric\nrepresentation,\n $$\n \\mbox{{\\goth\\char75}}(\\kappa)=F\\big({1\\over2},{1\\over2},1; \\kappa\\big)={1\\over\\pi\\kappa^{1\/4}}\\,\n Q_{-1\/2}\\big((\\kappa+1)\/2\\sqrt\\kappa\\big)\\,.\n $$\nA similar thing holds in the boson case of course.\n\\end{ignore}\n\\section{\\bf 8. Conclusion}\n\nIt does not seem possible to ellipticise integrals of the internal\nenergy, in particular the free energy (except for the circle), because\nthe modular properties are more complicated possessing nonzero cocycle\nfunctions. These are related to a known Lambert series connected with the\nEisenstein series and Zagier's expression for the period functions of\nthese is more neatly re--computed by a contour method. The use of the\nfully subtracted series $\\epsilon_t^{\\rm sub}$ allows a more compact\ntreatment.\n\nSufficient historical material has been exhibited to indicate that the\nSelberg--Chowla formula should be renamed the Epstein--Kober formula.\n\\newpage\n\\section{\\bf Appendix, Dirichlet Series.}\n\nIt is possible to give a fairly broad context to the inversion behaviour\n({\\it e.g. } modular properties), without too much restriction, in terms of\nDirichlet series. Although these ideas, by now classic, are associated\nwith Hecke, and have been expounded by Ogg [\\putref{Ogg,Ogg2}], for\nexample, the simplest case occurs earlier in other places. I mention\nKoshliakov, [\\putref{Kosh}], whose discussion I now paraphrase both for\nhistorical justice and interest.\n\nDefine two ordinary Dirichlet series (`zeta functions'),\n $$\n \\phi(s)=\\sum_{n=1}^\\infty {a_n\\over n^s}\\,,\\,\\,({\\rm Re\\,} s>\\nu_a)\\,,\\quad\n \\psi(s)=\\sum_{n=1}^\\infty {b_n\\over n^s}\\,,\\,\\,({\\rm Re\\,} s>\\nu_b>\\nu_a)\\,,\n $$\nfor some $\\nu$'s depending on the asymptotics of $a_n$ and $b_n$, and\n{\\it assume} that the $a_n$ and $b_n$ are such that the corresponding\n`cylinder kernels' satisfy the `formula of transformation',\n $$\n \\sum_{n=0}^\\infty a_n\\,e^{-nb\\rho}={a\\over\\rho^\\nu}\n \\sum_{n=0}^\\infty b_n\\,e^{-nb\/\\rho}\\,,\\quad \\nu=\\nu_a\\,,\\quad \\rho>0\\,,\n \\eql{trans2}$$\nfor fixed $a$ and $b$, where the numbers (not necessarily integers)\nof `zero modes' are taken to be,\n $$\n a_0=-\\phi(0)\\,\\quad b_0=-\\psi(0)\\,.\n $$\nThen Koshliakov, [\\putref{Kosh}], shows the equivalence of (\\puteqn{trans2})\nwith\n $$\n a\\,{\\Gamma(\\nu-s)\\,\\psi(\\nu-s)\\over b^{\\nu-s}}={\\Gamma(s)\\,\\phi(s)\\over b^s}\\,,\n \\eql{frel}$$\nby brute force using the summation function, $\\sigma(z)$, given by the\nBessel function expression,\n $$\n \\sigma(z)=-2ab\\,z^{(\\nu-1)\/2}\\,\\sum_{n=1}^\\infty {b_n\\over n^{(\\nu-1)\/2}}\\,\n K_{\\nu-1}\\big(2b\\sqrt{nz}\\big)\\,,\n $$\nand Heaviside's formula, (\\puteqn{hi}). $\\sigma(z)$ has the important property\nof possessing poles at $z=1,2,\\ldots$ with residues $a_1,a_2,\\ldots$ and\nalso has a cut along the negative $x$--axis with discontinuity\n$ab^\\nu2\\pi i\\psi(0)(-x)^{\\nu-1}\/\\Gamma(\\nu)$, $z=x+iy$.\n\nKoshliakov shows that $\\phi(s)$ is a single--valued holomorphic function\nhaving a first-order pole at $s=\\nu$,\n $$\n \\phi(s)\\sim-\\psi(0){ab^\\nu\\over\\Gamma(\\nu)}{1\\over s-\\nu}\\,,\n \\eql{phipole}$$\nand it is then a theorem that (\\puteqn{frel}) together with (\\puteqn{phipole}) is\nequivalent to (\\puteqn{trans2}), as can also be established easily, and\ninstructively, by Mellin transform. It follows, symmetrically, that\n$\\psi(s)$ is a single--valued holomorphic function having a first-order\npole at $s=\\nu$,\n $$\n \\psi(s)\\sim-\\phi(0){b^\\nu\\over a\\Gamma(\\nu)}{1\\over s-\\nu}\\,.\n $$\n\nKoshliakov also gives an integral form of $\\phi(s)$,\n $$\n \\phi(s)=-{\\sin\\pi s\\over\\pi}\\int_0^\\infty dx\\, {\\sigma(x)\\over x^s}\n \\,,\\quad {\\rm Re\\,} s<0\\,,\n \\eql{intphi}$$\nshowing that $\\phi(s)$ vanishes at negative integers in agreement with\n(\\puteqn{frel}).\n\nHecke, [\\putref{Hecke}], derives the same equivalence, but only for\n$a_n=b_n$, and eight years later. The more general case can also be found\nin the standard references Ogg [\\putref{Ogg,Ogg2}] and Weil,\n[\\putref{Weil2}], and no doubt elsewhere. A related formula involving\nRamanujan's `reciprocal functions' occurs in Hardy and Littlewood,\n[\\putref{HandL}].\n\nDefining cylinder--kernels, or theta--series, including the zero modes,\nby,\n $$\n \\Phi(\\beta)=\\sum_{n=0}^\\infty a_n\\,e^{-n\\beta}\\,,\\quad\n \\Psi(\\beta)=\\sum_{n=0}^\\infty b_n\\,e^{-n\\beta}\\,,\n $$\nthe reciprocal relation, (\\puteqn{trans2}), reads,\n $$\n \\Phi(b\\rho)={a\\over\\rho^\\nu}\\,\\Psi(b\/\\rho)\\,,\n \\eql{recip}\n $$\n(continued from the regions of convergence in $s$) and one has, in\nstandard fashion ({\\it cf } (\\puteqn{ell})),\n $$\n \\phi(s)={1\\over\\Gamma(s)}\\int_0^\\infty d\\beta\\,\\beta^{s-1}\\,\\big(\\Phi(\\beta)-a_0\\big)\\,,\n \\quad\n \\psi(s)={1\\over\\Gamma(s)}\\int_0^\\infty d\\beta\\,\\beta^{s-1}\n \\,\\big(\\Psi(\\beta)-b_0\\big)\\,.\n $$\n\nAs an example, the Dirichlet series corresponding to the classical\nEisenstein series is, from (\\puteqn{mell1}) or (\\puteqn{sig}),\n $$\n \\sum_{m=1}^\\infty \\sigma_{2t-1}(m)\\,m^{-s}=\\zeta_R(s)\\,\\zeta_R(s+1-2t)\\,,\n $$\nand this is a case discussed by Koshliakov, [\\putref{Kosh}] p.19, with\n$a_n=b_n=\\sigma_{2t-1}(n), \\nu=2t,a=(-1)^t,b=2\\pi,\\phi(0)=\\psi(0)=-B_{2t}\/4t$.\n\nKoshliakov also considers Jacobi elliptic cases that yield reciprocal\nidentities covered earlier by Glaisher [\\putref{Glaisher}].\n\n\\subsection{\\it Generalised Dirichlet series.}\nI have described above the classic Dirichlet set up. A certain, and\nsometimes only apparent, generalisation is obtained by replacing $n$ by\narbitrary sequences, $\\lambda_m$ and $\\mu_n$, and defining new $\\phi$ and\n$\\psi$ by\n $$\n \\phi(s)=\\sum_{m=1}^\\infty {a_m\\over \\lambda_m^s}\\,,\\quad\\quad\n \\psi(s)=\\sum_{n=1}^\\infty {b_n\\over \\mu_n^s}\\,,\n \\eql{newdir}$$\nwhich satisfy, {\\it by definition}, the functional equation\n(one of many possible),\n $$\n \\Gamma(\\delta-s)\\,\\psi(\\delta-s)=\\Gamma(s)\\,\\phi(s)\\,,\\quad \\delta\\in\\mbox{\\open\\char82}\\,,\n \\eql{frel2}$$\nwith corresponding heat--kernels,\n $$\n \\Phi(\\beta)=\\sum_{m=1}^\\infty a_m\\,e^{-\\lambda_m\\beta}\\,,\\quad\n \\Psi(\\beta)=\\sum_{n=1}^\\infty b_n\\,e^{-\\mu_n\\beta}\\,.\n \\eql{hkz}$$\n\nBochner, [\\putref{Bochner}], (see also Chandrasekharan and Narasimhan,\n[\\putref{CandN}], Knopp, [\\putref{Knopp2}] ), derives, by Mellin transforms,\nthe `modular relation',\n $$\n \\Phi(\\beta)-\\beta^{-\\delta}\\,\\Psi\\big({1\\over\\beta}\\big)=B(\\beta)\\,,\n \\eql{modrel}$$\nwhere the `residual function' $B(\\beta)$ is\n $$\n B(\\beta)={1\\over2\\pi i}\\int_{C({\\cal S})} ds\\,\\chi(s)\\,\\beta^{-s}=\\sum_{s\\in{\\cal S}}\\beta^{-s}\n {\\rm Res}\\,\\chi(s)\\,.\n $$\n$\\chi(s)$ is the joint continuation of the two sides of (\\puteqn{frel2}) and\nis assumed to have only simple poles for singularities confined to a\ncompact region, ${\\cal S}$, of the $s$--plane, usually on the real axis.\n$C({\\cal S})$ is a loop surrounding ${\\cal S}$. I especially note that if there\nexist higher--order poles then integer powers of $\\log \\beta$ occur.\n\nIn certain situations, {\\it e.g. } Koshliakov's, $B(\\beta)$ can be naturally\namalgamated with\nthe left--hand side to give an exact modular relation like (\\puteqn{trans2}),\n{\\it e.g. } Bochner,\n[\\putref{Bochner}] p.342. Typically one would include zero modes $\\lambda_0=0$\nand $\\mu_0=0$ in (\\puteqn{hkz}), but not in (\\puteqn{newdir}). (There seems to be\na misprint on p.342 of [\\putref{Bochner}] which has $\\lambda_0=\\mu_0=1$.)\n\nThere is another formula connected with this analysis that is useful. It\nis concerned with the representation of the {\\it modified} series,\n $$\n \\phi(s,w)=\\sum_{n=1}^\\infty {a_n\\over (\\lambda_n+w^2)^s}\\,\\quad{\\rm Re\\,} w>0.\n \\eql{asser}$$\n\nThere are many reasons why we should wish to add the number $w^2$. For\nexample it might correspond to a mass or simply be added for extra\nflexibility as when calculating heat--kernel expansion coefficients via\nresolvents or it might be another part of the denominator as in Epstein's\nderivation of (\\puteqn{kob1}); see [\\putref{DandKi}].\n\nThe representation alluded to is, (Berndt [\\putref{Berndt4}]),\n $$\n \\Gamma(s)\\,\\phi(s,w)=R(s,w)+2\\sum_{n=1}^\\infty b_n\\,\\bigg({\\mu_n\\over w^2}\n \\bigg)^{(s-\\delta)\/2}\\,K_{s-\\delta}\\big(2w\\sqrt\\mu_n\\big)\\,,\n \\eql{berndt}$$\nwhere $s$ is such that the summation converges absolutely and $R(s,w)$ is\ngiven by\n $$\n R(s,w)=\\sum_{s'\\in\\,{\\cal S}} \\Gamma(s-s')\\,w^{2s'-2s}{\\rm Res}\\,\\, \\chi(s')\\,.\n \\eql{resi}$$\n\nBerndt's second proof of (\\puteqn{berndt}) proceeds via the\ncylinder--kernels, $\\Phi$ and $\\Psi$. I prefer his first, Mellin, proof,\n[\\putref{Berndt6}]. Of course, the information employed is the same. I give\nthis proof out of interest.\n\nThe Mellin transform relation between $\\phi(s)$ and the $\\phi(s,w)$ of\n(\\puteqn{asser}), is\n $$\n \\Gamma(s)\\,\\phi(s,w)={1\\over2\\pi i}\\int_{c-i\\infty}^{c+i\\infty}ds'\\,w^{2s'-2s}\n \\Gamma(s-s')\\Gamma(s')\\,\\phi(s')\n $$\nwhere $c$ is chosen sufficiently positive to make both $\\phi(s)$ and\n$\\psi(s)$ converge. The vertical contour is now moved to the left so as\nto run from $\\delta-c-i\\infty$ to $\\delta-c+i\\infty$ picking up the residues,\n$R(s,w)$, (\\puteqn{resi}), on the way. The region, ${\\cal S}$, lies between\nthese vertical lines and also the horizontal pieces contribute zero.\n\nOn the shifted vertical contour make the coordinate change $s'\\to\\delta-s'$\nand use (\\puteqn{frel2}) {\\it i.e. } $\\chi(\\delta-s')=\\Gamma(s')\\psi(s')$, valid in this\nrange, and so for this part I find,\n $$\\eqalign{\n \\sum_n b_n {1\\over2\\pi i}&\\int_{c-i\\infty}^{c+i\\infty}ds'\\,\\Gamma(s'-\\delta+s)\n \\Gamma(s')\\mu_n^{-s'} w^{2\\delta-2s'-2s}\\cr\n&=2\\sum_{n=1}^\\infty b_n\\,\\bigg({\\mu_n\\over w^2}\n \\bigg)^{(s-\\delta)\/2}\\,K_{s-\\delta}\\big(2w\\sqrt\\mu_n\\big)\\,,}\n $$\non using (\\puteqn{hi}). Hence (\\puteqn{berndt}) has been established.\n\nIn specific situations (\\puteqn{berndt}) is quite familiar in the $\\zeta$--function\\ area\nand in background and finite--temperature quantum field theory.\n\nAn illustrative example is the simplest, diagonal Epstein function,\n $$\n \\phi(s)=Z_p(s)\\equiv\\sumdasht{{\\bf m}=-\\infty}{\\infty}\n {1\\over({\\bf m.m})^s}\\,,\\quad {\\bf m}=(m_1,m_2,\\ldots,m_p)\\,,\n \\eql{diagep}$$\nwhich obeys the functional equation (Epstein [\\putref{Epstein}]),\n $$\n \\Gamma(s)\\,Z_p(s)=\\pi^{2s-p\/2}\\,\\Gamma(p\/2-s)\\,Z_p(p\/2-s)\\,,\n \\eql{epfunc}$$\nso that, see (\\puteqn{frel2}), $\\delta=p\/2$, $a_n=1$, $\\lambda_m = {\\bf m.m}$,\n$b_n=\\pi^{p\/2}$, $\\mu_n=\\pi^2\\,{\\bf m.m}$.\n\n$\\Gamma(s)\\,Z_p(s)$ has poles at $s=0$ and $s=p\/2$ and (\\puteqn{berndt}) gives\nfor the modified, or inhomogeneous, or `massive' Epstein function,\n$Z_p(s,w)$,\n $$\\eqalign{\n \\sumdasht{{\\bf m}=-\\infty}{\\infty}\n {1\\over({\\bf m.m}+w^2)^s}=-w^{-2s}&+{\\Gamma(p\/2)\\Gamma(s-p\/2)\\over\\Gamma(s)}\n \\,w^{p-2s}\\cr\n & +{2\\pi^s\\over\\Gamma(s)}\\sumdasht{{\\bf m}=-\\infty}{\\infty}\n \\bigg({|{\\bf m}|\\over w}\\bigg)^{s-p\/2}\\,\n K_{s-p\/2}\\big(2w\\pi|{\\bf m}|\\big)\\,,}\n \\eql{modep}$$\nwhere the dash means to omit the ${\\bf m=0}$ term. Note that the $w^{-2s}$ could be\nincluded in the left--hand sum if this were extended to include ${\\bf m=0}$.\n\nThe right--hand side can be taken as the continuation of the left--hand\nside showing the single pole at $s=p\/2$, as is correct and may be proved in\nother ways.\n\nActually (\\puteqn{modep}) is shown much more directly using the Jacobi\ninversion relation, which is, of course, how Epstein obtained\n(\\puteqn{epfunc}). Equivalent to Jacobi inversion is Poisson summation.\n\nEquations like (\\puteqn{modep}) also occur in lattice summations. Some references,\nbut none earlier than those already quoted, are given in the review by Glasser\nand Zucker [\\putref{GandZ}] p.109.\n\nA similar formula holds for the general Epstein function which may, or may\nnot, have a pole at $s=p\/2$ and may, or may not, vanish at $s=0$. The works\n[\\putref{Kirsten1,Elizalde2,Elizalde3}], for example, can be consulted for details,\nsome further references and physical applications.\n\nWhen $p=1$, $\\phi(s)=2\\zeta_R(2s)$ and I recover an earlier, well--known\nformula, [\\putref{Watson,Kober}], which yields, after putting $w=u\\,m_2$\nand summing over $m_2$, the formula for the Epstein $\\zeta$--function\\ mentioned\nearlier; see (\\puteqn{kob1}).\n\nIt is left as an exercise to show that the non-diagonal formula,\n(\\puteqn{kob1}), can be obtained in the same way using reciprocal relations\nfor the Hurwitz $\\zeta$--function, or, equivalently, the Lipshitz formula ({\\it cf } Epstein,\n[\\putref{Epstein}], Kober [\\putref{Kober}] p.622).\n\nAs an application of (\\puteqn{modep}) it might be helpful if I present a\nstandard derivation of the usual statistical free energy mode sum from\nthe thermal $\\zeta$--function\\ expression, for a reasonably general system. The\nstarting point is the determinant form,\n $$\n F={i\\over2}\\lim_{s\\to0}{\\zeta(s,\\beta)\\over s}\n =-{1\\over2\\beta}\\lim_{s\\to0}{1\\over s}\n \\sumdasht{{m=-\\infty\\atop n}}\n{\\infty}{d_n\\over\n \\big(\\omega_n^2+4\\pi^2m^2\/\\beta^2\\big)^s}\\,,\n \\eql{eff2}$$\nwhere $\\omega_n^2$ and $d_n$ are the eigenvalues and degeneracies of the\nappropriate operator (related to the Laplacian) on the spatial section,\n${\\cal M}$. The $\\omega_n$ are the single--particle energies. The dash here means\nthat the denominator should never be zero ({\\it cf } (\\puteqn{eff})) but the sum\nincludes $m=0$. For simplicity it is assumed that there is no spatial zero\nmode and so the dash can be removed.\n\nNow consider (\\puteqn{modep}) for $p=1$. Include the $w^{-2s}$ with the\nleft--hand sum, set $w=\\beta\\omega_n\/2\\pi$, multiply by the degeneracy, $d_n$\nand sum over $n$. The limit (\\puteqn{eff2}) is then easily taken since\n$s\\Gamma(s)=\\Gamma(s+1)$ and only the $K_{1\/2}$ Bessel appears. There is\nnothing deep about this calculation and the result is a standard form,\n $$\\eqalign{\n F=&{1\\over2}\\,\\zeta_{\\cal M}(-1\/2)+{1\\over\\beta}\\sum_n d_n\n \\sum_{m=1}^\\infty{e^{-m\\omega_n\\beta}\\over m}\\cr\n =&{1\\over2}\\,\\zeta_{\\cal M}(-1\/2)+{1\\over\\beta}\\sum_n d_n\\log(1-e^{-\\omega_n\\beta})\\,.}\n \\eql{modesum}$$\nTo avoid specifically field theoretic problems, it has been assumed that\n$\\zeta_{\\cal M}(-1\/2)$ exists. This is true for the cases discussed in this\npaper.\n\nThe result, (\\puteqn{modesum}), is an analogue of the (first) Kronecker\nlimit--formula applied to the generalised torus S$^1\\times{\\cal M}$; see\n[\\putref{DandA}]. The original formula relates to ${\\cal M}=$S$^1$ and, in this\ncase, there exists a functional relation, (\\puteqn{frel2}), so that the\nlimit--formula is often expressed in terms of the remainder about the\npole in $\\zeta(s,\\beta)$ at $s=1$. In the general case this is not possible.\n\nFinally, if one introduces the standard arithmetic quantity $r_p(n)$,\nequal to the number of representations of the integer $n$ as the sum of\n$p$ squares, then the diagonal Epstein function, (\\puteqn{diagep}), is\nobtained by writing,\n $$\n \\phi(s)=Z_p(s)=\\sum_n{r_{p}(n)\\over n^s}\\,,\n $$\nand is treated as such by Koshliakov who also gives the `exact', product\nforms of $Z_p(s)$ for $p=2$ and $p=4$.\n\n\\begin{ignore}\nThe other useful formula is a generalisation of one occurring in the\ntheory of Riesz means which has already been encountered, essentially in\n(\\puteqn{mint}) and (\\puteqn{soln}).\n\nWith the same notation as above, one has, for $h\\ge0$,\n $$\n {1\\over\\Gamma(h+1)}\\sum_{\\lambda_n}a_n\\,\\theta(x-\\lambda_n)\\,(x-\\lambda_n)^h\n ={1\\over2\\pi i}\\int_{c-i\\infty}\n ^{c+i\\infty}ds\\,{\\Gamma(s)\\,\\phi(s)\\,x^{s+h}\\over\\Gamma(h+1+s)}\\,,\n $$\nwhere $\\theta$ is Heaviside's step function.\n\\end{ignore}\n\\newpage\n\\section{\\bf References.}\n\n\\begin{putreferences}\n \\reference{RandA}{Rao,M.B. and Ayyar,M.V. \\jims{15}{1923\/24}{150}.}\n \\reference{DandK2}{Dowker,J.S. and Kirsten,K. {\\it Elliptic aspects of\n statistical mechanics on spheres} ArXiv:0807.1995.}\n \\reference{DandK1}{Dowker,J.S. and Kirsten,K. \\np {638}{2002}{405}.}\n \\reference{GandS}{Gel'fand, I.M. and Shilov, G.E. {\\it Generalised Functions}, vol.1\n (Academic Press, New York, 1964). }\n \\reference{DandA}{Dowker,J.S. and Apps, J.S. \\cqg{12}{1995}{1363}.}\n \\reference{Weil}{Weil,A., {\\it Elliptic functions according to Eisenstein and\n Kronecker}, Springer, Berlin, 1976.}\n \\reference{Ling}{Ling,C-H. {\\it SIAM J.Math.Anal.} {\\bf5} (1974) 551.}\n \\reference{Ling2}{Ling,C-H. {\\it J.Math.Anal.Appl.}(1988).}\n \\reference{BMO}{Brevik,I., Milton,K.A. and Odintsov, S.D. {\\it Entropy bounds in\n $R\\times S^3$ geometries}. hep-th\/0202048.}\n \\reference{KandL}{Kutasov,D. and Larsen,F. {\\it JHEP} 0101 (2001) 1.}\n \\reference{KPS}{Klemm,D., Petkou,A.C. and Siopsis {\\it Entropy\n bounds, monoticity properties and scaling in CFT's}. hep-th\/0101076.}\n \\reference{DandC}{Dowker,J.S. and Critchley,R. \\prD{15}{1976}{1484}.}\n \\reference{AandD}{Al'taie, M.B. and Dowker, J.S. \\prD{18}{1978}{3557}.}\n \\reference{Dow1}{Dowker,J.S. \\prD{37}{1988}{558}.}\n \\reference{Dow3}{Dowker,J.S.\\prD{28}{1983}{3013}.}\n \\reference{DandK}{Dowker,J.S. and Kennedy,G. \\jpa{11}{1978}{895}.}\n \\reference{Dow2}{Dowker,J.S. \\cqg{1}{1984}{359}.}\n \\reference{DandKi}{Dowker,J.S. and Kirsten, K. {\\it Comm. in Anal. and Geom.\n }{\\bf7}(1999) 641.}\n \\reference{DandKe}{Dowker,J.S. and Kennedy,G.\n \\jpa{11}{1978}{895}.} \\reference{Gibbons}{Gibbons,G.W. \\pl{60A}{1977}{385}.}\n \\reference{Cardy}{Cardy,J.L. \\np{366}{1991}{403}.}\n \\reference{ChandD}{Chang,P. and\n Dowker,J.S. \\np{395}{1993}{407}.}\n \\reference{DandC2}{Dowker,J.S. and\n Critchley,R. \\prD{13}{1976}{224}.}\n \\reference{Camporesi}{Camporesi,R.\n \\prp{196}{1990}{1}.}\n \\reference{BandM}{Brown,L.S. and Maclay,G.J.\n \\pr{184}{1969}{1272}.}\n \\reference{CandD}{Candelas,P. and Dowker,J.S.\n \\prD{19}{1979}{2902}.}\n \\reference{Unwin1}{Unwin,S.D. Thesis. University of\n Manchester. 1979.} \\reference{Unwin2}{Unwin,S.D. \\jpa{13}{1980}{313}.}\n \\reference{DandB}{Dowker,J.S. and Banach,R. \\jpa{11}{1979}{}.}\n \\reference{Obhukov}{Obhukov,Yu.N. \\pl{109B}{1982}{195}.}\n \\reference{Kennedy}{Kennedy,G. \\prD{23}{1981}{2884}.}\n \\reference{CandT}{Copeland,E.\n and Toms,D.J. \\np {255}{1985}{201}.} \\reference{ELV}{Elizalde,E., Lygren, M.\n and Vassilevich, D.V. \\jmp {37}{1996}{3105}.}\n \\reference{Malurkar}{Malurkar,S.L. {\\it J.Ind.Math.Soc} {\\bf16} (1925\/26) 130.}\n \\reference{Glaisher}{Glaisher,J.W.L. {\\it Messenger of Math.} {\\bf18} (1889)\n 1.}\n \\reference{Anderson}{Anderson,A. \\prD{37}{1988}{536}.}\n \\reference{CandA}{Cappelli,A. and D'Appollonio,\\pl{487B}{2000}{87}.}\n \\reference{Wot}{Wotzasek,C. \\jpa{23}{1990}{1627}.}\n \\reference{RandT}{Ravndal,F. and Tollesen,D. \\prD{40}{1989}{4191}.}\n \\reference{SandT}{Santos,F.C. and Tort,A.C. \\pl{482B}{2000}{323}.}\n \\reference{FandO}{Fukushima,K. and Ohta,K. {\\it Physica} {\\bf A299} (2001) 455.}\n \\reference{GandP}{Gibbons,G.W. and Perry,M. \\prs{358}{1978}{467}.}\n \\reference{Dow4}{Dowker,J.S. {\\it Zero modes, entropy bounds and partition\n functions.} hep-th\\break \/0203026.}\n \\reference{Rad}{Rademacher,H. {\\it Topics in analytic number theory,}\n (Springer- Verlag, Berlin, 1973).}\n \\reference{Halphen}{Halphen,G.-H. {\\it Trait\\'e des Fonctions Elliptiques}, Vol 1,\n Gauthier-Villars, Paris, 1886.}\n \\reference{CandW}{Cahn,R.S. and Wolf,J.A. {\\it Comm.Mat.Helv.} {\\bf 51} (1976) 1.}\n \\reference{Berndt}{Berndt,B.C. \\rmjm{7}{1977}{147}.}\n \\reference{Hurwitz}{Hurwitz,A. \\ma{18}{1881}{528}.}\n \\reference{Hurwitz2}{Hurwitz,A. {\\it Mathematische Werke} Vol.I. Basel,\n Birkhauser, 1932.}\n \\reference{Berndt2}{Berndt,B.C. \\jram{303\/304}{1978}{332}.}\n \\reference{RandA}{Rao,M.B. and Ayyar,M.V. \\jims{15}{1923\/24}{150}.}\n \\reference{Hardy}{Hardy,G.H. \\jlms{3}{1928}{238}.}\n \\reference{TandM}{Tannery,J. and Molk,J. {\\it Fonctions Elliptiques},\n Gauthier-Villars, Paris, 1893--1902.}\n \\reference{schwarz}{Schwarz,H.-A. {\\it Formeln und Lehrs\\\"atzen zum Gebrauche..},\n Springer 1893.(The first edition was 1885.) The French translation by\n Henri Pad\\'e is, Gauthier-Villars, Paris, 1894}\n \\reference{Hancock}{Hancock,H. {\\it Theory of elliptic functions}, Vol I.\n Wiley, New York 1910.}\n \\reference{watson}{Watson,G.N. \\jlms{3}{1928}{216}.}\n \\reference{MandO}{Magnus,W. and Oberhettinger,F. {\\it Formeln und S\\\"atze},\n Springer-Verlag, Berlin 1948.}\n \\reference{Klein}{Klein,F. {\\it }.}\n \\reference{AandL}{Appell,P. and Lacour,E. {\\it Fonctions Elliptiques},\n Gauthier-Villars,\n Paris, 1897.}\n \\reference{HandC}{Hurwitz,A. and Courant,C. {\\it Allgemeine Funktionentheorie},\n Springer,\n Berlin, 1922.}\n \\reference{WandW}{Whittaker,E.T. and Watson,G.N. {\\it Modern analysis},\n Cambridge 1927.}\n \\reference{SandC}{Selberg,A. and Chowla,S. \\jram{227}{1967}{86}. }\n \\reference{zucker}{Zucker,I.J. {\\it Math.Proc.Camb.Phil.Soc} {\\bf 82 }(1977) 111.}\n \\reference{glasser}{Glasser,M.L. {\\it Maths.of Comp.} {\\bf 25} (1971) 533.}\n \\reference{GandW}{Glasser, M.L. and Wood,V.E. {\\it Maths of Comp.} {\\bf 25} (1971)\n 535.}\n \\reference{greenhill}{Greenhill,A,G. {\\it The Applications of Elliptic\n Functions}, MacMillan, London, 1892.}\n \\reference{Weierstrass}{Weierstrass,K. {\\it J.f.Mathematik (Crelle)}\n {\\bf 52} (1856) 346.}\n \\reference{Weierstrass2}{Weierstrass,K. {\\it Mathematische Werke} Vol.I,p.1,\n Mayer u. M\\\"uller, Berlin, 1894.}\n \\reference{Fricke}{Fricke,R. {\\it Die Elliptische Funktionen und Ihre Anwendungen},\n Teubner, Leipzig. 1915, 1922.}\n \\reference{Konig}{K\\\"onigsberger,L. {\\it Vorlesungen \\\"uber die Theorie der\n Elliptischen Funktionen}, \\break Teubner, Leipzig, 1874.}\n \\reference{Milne}{Milne,S.C. {\\it The Ramanujan Journal} {\\bf 6} (2002) 7-149.}\n \\reference{Schlomilch}{Schl\\\"omilch,O. {\\it Ber. Verh. K. Sachs. Gesell. Wiss.\n Leipzig} {\\bf 29} (1877) 101-105; {\\it Compendium der h\\\"oheren Analysis},\n Bd.II, 3rd Edn, Vieweg, Brunswick, 1878.}\n \\reference{BandB}{Briot,C. and Bouquet,C. {\\it Th\\`eorie des Fonctions Elliptiques},\n Gauthier-Villars, Paris, 1875.}\n \\reference{Dumont}{Dumont,D. \\aim {41}{1981}{1}.}\n \\reference{Andre}{Andr\\'e,D. {\\it Ann.\\'Ecole Normale Superior} {\\bf 6} (1877) 265;\n {\\it J.Math.Pures et Appl.} {\\bf 5} (1878) 31.}\n \\reference{Raman}{Ramanujan,S. {\\it Trans.Camb.Phil.Soc.} {\\bf 22} (1916) 159;\n {\\it Collected Papers}, Cambridge, 1927}\n \\reference{Weber}{Weber,H.M. {\\it Lehrbuch der Algebra} Bd.III, Vieweg,\n Brunswick 190 3.}\n \\reference{Weber2}{Weber,H.M. {\\it Elliptische Funktionen und algebraische Zahlen},\n Vieweg, Brunswick 1891.}\n \\reference{ZandR}{Zucker,I.J. and Robertson,M.M.\n {\\it Math.Proc.Camb.Phil.Soc} {\\bf 95 }(1984) 5.}\n \\reference{JandZ1}{Joyce,G.S. and Zucker,I.J.\n {\\it Math.Proc.Camb.Phil.Soc} {\\bf 109 }(1991) 257.}\n \\reference{JandZ2}{Zucker,I.J. and Joyce.G.S.\n {\\it Math.Proc.Camb.Phil.Soc} {\\bf 131 }(2001) 309.}\n \\reference{zucker2}{Zucker,I.J. {\\it SIAM J.Math.Anal.} {\\bf 10} (1979) 192,}\n \\reference{BandZ}{Borwein,J.M. and Zucker,I.J. {\\it IMA J.Math.Anal.} {\\bf 12}\n (1992) 519.}\n \\reference{Cox}{Cox,D.A. {\\it Primes of the form $x^2+n\\,y^2$}, Wiley, New York,\n 1989.}\n \\reference{BandCh}{Berndt,B.C. and Chan,H.H. {\\it Mathematika} {\\bf42} (1995) 278.}\n \\reference{EandT}{Elizalde,R. and Tort.hep-th\/}\n \\reference{KandS}{Kiyek,K. and Schmidt,H. {\\it Arch.Math.} {\\bf 18} (1967) 438.}\n \\reference{Oshima}{Oshima,K. \\prD{46}{1992}{4765}.}\n \\reference{greenhill2}{Greenhill,A.G. \\plms{19} {1888} {301}.}\n \\reference{Russell}{Russell,R. \\plms{19} {1888} {91}.}\n \\reference{BandB}{Borwein,J.M. and Borwein,P.B. {\\it Pi and the AGM}, Wiley,\n New York, 1998.}\n \\reference{Resnikoff}{Resnikoff,H.L. \\tams{124}{1966}{334}.}\n \\reference{vandp}{Van der Pol, B. {\\it Indag.Math.} {\\bf18} (1951) 261,272.}\n \\reference{Rankin}{Rankin,R.A. {\\it Modular forms} (CUP, Cambridge, 1977).}\n \\reference{Rankin2}{Rankin,R.A. {\\it Proc. Roy.Soc. Edin.} {\\bf76 A} (1976) 107.}\n \\reference{Skoruppa}{Skoruppa,N-P. {\\it J.of Number Th.} {\\bf43} (1993) 68 .}\n \\reference{Down}{Dowker.J.S. \\np {104}{2002}{153}.}\n \\reference{Eichler}{Eichler,M. \\mz {67}{1957}{267}.}\n \\reference{Zagier}{Zagier,D. \\invm{104}{1991}{449}.}\n \\reference{KandZ}{Kohnen, W. and Zagier,D. {\\it Modular forms with rational periods.}\n in {\\it Modular Forms} ed. by Rankin, R.A. (Ellis Horwood, Chichester, 1984).}\n \\reference{Lang}{Lang,S. {\\it Modular Forms}, (Springer, Berlin, 1976).}\n \\reference{Kosh}{Koshliakov,N.S. {\\it Mess.of Math.} {\\bf 58} (1928) 1.}\n \\reference{BandH}{Bodendiek, R. and Halbritter,U. \\amsh{38}{1972}{147}.}\n \\reference{Smart}{Smart,L.R., \\pgma{14}{1973}{1}.}\n \\reference{Grosswald}{Grosswald,E. {\\it Acta. Arith.} {\\bf 21} (1972) 25.}\n \\reference{Kata}{Katayama,K. {\\it Acta Arith.} {\\bf 22} (1973) 149.}\n \\reference{Ogg}{Ogg,A. {\\it Modular forms and Dirichlet series} (Benjamin, New York,\n 1969).}\n \\reference{Bol}{Bol,G. \\amsh{16}{1949}{1}.}\n \\reference{Epstein}{Epstein,P. \\ma{56}{1902}{615}.}\n \\reference{Petersson}{Petersson, H. \\ma {116}{1939}{401}.}\n \\reference{Serre}{Serre,J-P. {\\it A Course in Arithmetic}, Springer, New York,\n 1973.}\n \\reference{Schoenberg}{Schoenberg,B., {\\it Elliptic Modular Functions},\n Springer, Berlin, 1974.}\n \\reference{Apostol}{Apostol,T.M. \\dmj {17}{1950}{147}.}\n \\reference{Ogg2}{Ogg,A. {\\it Lecture Notes in Math.} {\\bf 320} (1973) 1.}\n \\reference{Knopp}{Knopp,M.I. \\dmj {45} {1978}{47}.}\n \\reference{Knopp2}{Knopp,M.I. \\invm {}{1994}{361}.}\n \\reference{Knopp3}{Knopp,M.I. {\\it Lecture Notes in Mathematics} {\\bf 1383}(1989) 111 .}\n \\reference{LandZ}{Lewis,J. and Zagier,D. \\aom{153}{2001}{191}.}\n \\reference{HandK}{Husseini, S.Y. and Knopp, M.I. {\\it Illinois.J.Math.} {\\bf 15} (1971)\n 565.}\n \\reference{Kober}{Kober,H. \\mz{39}{1934-5}{609}.}\n \\reference{HandL}{Hardy,G.H. and Littlewood, \\am{41}{1917}{119}.}\n \\reference{Watson}{Watson,G.N. \\qjm{2}{1931}{300}.}\n \\reference{SandC2}{Chowla,S. and Selberg,A. {\\it Proc.Nat.Acad.} {\\bf 35}\n (1949) 371.}\n \\reference{Landau}{Landau, E. {\\it Lehre von der Verteilung der Primzahlen},\n (Teubner, Leipzig, 1909).}\n \\reference{Berndt4}{Berndt,B.C. \\tams {146}{1969}{323}.}\n \\reference{Berndt3}{Berndt,B.C. \\tams {}{}{}.}\n \\reference{Bochner}{Bochner,S. \\aom{53}{1951}{332}.}\n \\reference{Weil2}{Weil,A.\\ma{168}{1967}{}.}\n \\reference{CandN}{Chandrasekharan,K. and Narasimhan,R. \\aom{74}{1961}{1}.}\n \\reference{Rankin3}{Rankin,R.A. {\\it Proc.Glas.Math.Assoc.} {\\bf 1} (1953) 149.}\n \\reference{Berndt6}{Berndt,B.C. {\\it Trans.Edin.Math.Soc} {\\bf 10} (1967) 309.}\n \\reference{Elizalde}{Elizalde,E. {\\it Ten Physical Applications of Spectral\n Zeta Function Theory}, \\break (Springer, Berlin, 1995).}\n \\reference{Allen}{Allen,B., Folacci,A. and Gibbons,G.W. \\pl{189}{1987}{304}.}\n \\reference{Krazer}{Krazer,A. {\\it Lehrbuch der Thetafunktionen} (Teubner,\n Leipzig, 1903)}\n \\reference{Elizalde3}{Elizalde,E. {\\it J.Comp.and Appl. Math.} {\\bf 118} (2000) 125.}\n \\reference{Elizalde2}{Elizalde,E., Odintsov.S.D, Romeo, A. and Bytsenko, A.A and\n Zerbini,S.\n {\\it Zeta function regularisation}, (World Scientific, Singapore, 1994).}\n \\reference{Eisenstein}{Eisenstein,B. \\jram {35}{1847}{153}.}\n \\reference{Hecke}{Hecke,E. \\ma{112}{1936}{664}.}\n \\reference{Terras}{Terras,A. {\\it Harmonic analysis on Symmetric Spaces} (Springer,\n New York, 1985).}\n \\reference{BandG}{Bateman,P.T. and Grosswald,E. {\\it Acta Arith.} {\\bf 9} (1964) 365.}\n \\reference{Deuring}{Deuring,M. \\aom{38}{1937}{585}.}\n \\reference{Guinand2}{Guinand, A.P. \\qjm{6}{1955}{156}.}\n \\reference{Guinand}{Guinand, A.P. \\qjm {15}{1944}{11}.}\n \\reference{Minak}{Minakshisundaram.}\n \\reference{Mordell}{Mordell,J. \\qjm {1}{1930}{77}.}\n \\reference{GandZ}{Glasser,M.L. and Zucker,I.J. ``Lattice Sums\" in Theoretical\n Chemistry, Advances and Perspectives, vol.5, ed. D.Henderson, New York,\n 1986, 67 .}\n \\reference{Landau2}{Landau,E. \\jram{}{1903}{64}.}\n \\reference{Kirsten1}{Kirsten,K. \\jmp{35}{1994}{459}.}\n \\reference{Sommer}{Sommer,J. {\\it Vorlesungen \\\"uber Zahlentheorie} (1907,Teubner,Leipzig).\n French edition 1913 .}\n \\reference{Reid}{Reid,L.W. {\\it Theory of Algebraic Numbers}, (1910,MacMillan,New York).}\n\\end{putreferences}\n\\bye\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\\label{sec:intro}\n\nLet $n \\geq 1$ be an integer and $\\xi \\in \\lR$.\nWe denote by $w_n(\\xi)$ (resp.\\ $w_n ^ {*}(\\xi)$) the supremum of the real numbers $w$ (resp.\\ $w ^ {*}$) which satisfy\n\\begin{equation*}\n\t0 < \\abs{P(\\xi)} \\leq H(P)^{- w} \\quad (\\text{resp.\\ } 0 < \\abs{\\xi - \\alpha} \\leq H(\\alpha)^{- w ^ {*} - 1})\n\\end{equation*}\nfor infinitely many integer polynomials $P(X)$ of degree at most $n$ (resp.\\ algebraic numbers $\\alpha \\in \\lC$ of degree at most $n$).\nHere, $H(P)$ is defined to be the maximum of the absolute values of the coefficients of $P(X)$ and $H(\\alpha)$ is equal to $H(Q)$, where $Q(X)$ is the minimal polynomial of $\\alpha$ over $\\lZ$.\nThe functions $w_n$ and $w_n ^ {*}$ are called \\textit{Diophatine exponents}.\nLiouville proved that $w_1 (\\alpha) \\leq \\deg \\alpha - 1$ for any algebraic real number $\\alpha$.\nRoth \\cite{Roth55} improved Liouville Theorem, that is, he showed that $w_1 (\\alpha) = 1$ for all algebraic irrational real numbers $\\alpha$.\nFor higher Diophantine exponents, from the Schmidt Subspace Theorem, it is known that\n\\begin{equation*}\n\tw_n (\\alpha ) = w_n ^ {*} (\\alpha ) = \\min \\{ n, d-1 \\},\n\\end{equation*}\nwhere $n \\geq 1$ is an integer and $\\alpha$ is an algebraic real number of degree $d$ (see \\cite[Section 3]{Bugeaud04}).\n\nLet $p$ be a prime, $q=p^s$, where $s \\geq 1$ is an integer and $\\lF_q$ be the finite field containing $q$ elements.\nWe denote by $\\lF_q [T], \\lF_q (T), $ and $\\lF_q ((T ^ {- 1}))$ respectively the ring of polynomials, the field of rational functions, the field of Laurent series in $T ^ {- 1}$ over $\\lF_q$.\nWe consider analogues of Diophantine exponents $w_n$ and $w_n ^ {*}$ for Laurent series (see Section \\ref{sec:notation} for the precisely definition).\nAs an analogue of Liouville Theorem, Mahler \\cite{Mahler49} showed $w_1 (\\alpha) \\leq \\deg \\alpha - 1$ for any algebraic Laurent series $\\alpha \\in \\lF_q((T^{-1}))$.\nHowever, there exist counter examples of analogue of Roth Theorem for Laurent series over a finite field.\nLet $r_1 = p ^ {t_1}$, where $t_1 \\geq 1$ is an integer.\nHe proved that the Laurent series $\\alpha_1 := \\sum_{n = 0}^{\\infty}T ^ {- r_1 ^ n}$ is algebraic of degree $r_1$ with $w_1 (\\alpha_1) = r_1 - 1$.\nAfter that many authors investigated rational approximations (e.g.\\ \\cite{Firicel13, Mathan92, Schmidt00, Thakur99}) and algebraic approximations (e.g.\\ \\cite{Ooto18, Thakur11, Thakur13}) for algebraic Laurent series.\nIn particular, the second author (\\cite{Ooto18}) proved that, for any rational number $w > 2n-1$, there exists an algebraic Laurent series $\\alpha \\in \\lF_q((T^{-1}))$ such that\n\\begin{equation*}\n\tw_1(\\alpha) = w_1^{*}(\\alpha) = \\ldots = w_n(\\alpha) = w_n^{*}(\\alpha) = w.\n\\end{equation*}\nThe first purpose of this paper is to consider Problem 2.2 in \\cite{Ooto18}.\n\n\n\n\\begin{prob}\\label{prob:main1}\n\tIs it true that\n\t\\begin{equation*}\n\t\tw_n(\\alpha) = w_n^{*}(\\alpha)\n\t\\end{equation*}\n\tfor an integer $n \\geq 1$ and an algebraic Laurent series $\\alpha \\in \\lF_q((T^{-1}))$?\n\\end{prob}\n\n\n\nNote that it is well-known and easy to see that $w_1(\\xi) = w_1^{*}(\\xi)$ for all $\\xi \\in \\lF_q((T^{-1}))$.\nWe give a negative answer of Problem \\ref{prob:main1} for any $n \\geq 2$ and $q \\geq 4$.\n\n\n\n\\begin{thm}\\label{thm:neq_alg}\n\tLet $n \\geq 2$ be an integer, $q = p^s$, and $r = q^t$ with $s, t \\geq 1$ be integers and \n\t\\begin{equation*}\n\t\tr \\geq \\frac{3n + 2 + \\sqrt{9n^2 + 4n + 4}}{2}.\n\t\\end{equation*}\n\tIf $q \\geq 4$, then there exist algebraic Laurent series $\\alpha \\in \\lF_q ((T ^ {- 1}))$ with $\\deg \\alpha = r+1$ such that\n\t\\begin{equation*}\n\t\tw_n (\\alpha) \\neq w_n ^ {*} (\\alpha)\n\t\\end{equation*}\n\\end{thm}\n\n\n\nLet $\\alpha$ be in $\\lF_q ((T ^ {- 1}))$ and $r = p ^ t$, where $t \\geq 0$ is an integer.\nWe call $\\alpha$ \\textit{hyperquadratic} if $\\alpha$ is irrational and there exists $(A, B, C, D) \\in \\lF_q [T] ^ 4 \\setminus \\{ {\\bf 0} \\}$ such that\n\\begin{equation}\\label{eq:hyperquad}\n\tA \\alpha^{r+1} + B \\alpha ^ r + C \\alpha + D = 0.\n\\end{equation}\nFor example, quadratic Laurent series are hyperquadratic, $\\alpha_1$ is hyperquadratic and satisfies $T \\alpha_1^r - T \\alpha_1 + 1 = 0$.\nBaum and Sweet \\cite{Baum76} started a study of hyperquadratic continued fractions in characteristic two.\nIn 1986, Mills and Robbins \\cite{Mills86} showed the existence of hyperquadratic continued fractions with all partial quotient of degree one in odd characteristic with a prime field.\nWe refer the reader to \\cite{Lasjaunias172} for a survey of recent works of hyperquadratic continued fractions.\nWe observe that the degree of hyperquadratic Laurent series which satisfies \\eqref{eq:hyperquad} is less than or equal to $r+1$.\nHowever, the exact degree of these is open.\nIn this paper, we discuss the following problem.\n\n\n\n\\begin{prob}\n\tGiven a hyperquadratic Laurent series $\\alpha \\in \\lF_q((T^{-1}))$, determine the exact degree of $\\alpha$.\n\\end{prob}\n\n\n\nIn order to determine the exact degree, we approach to quadratic approximations.\nSince $w_2(\\alpha) \\leq \\deg \\alpha - 1$ for $\\alpha$ as in \\eqref{eq:hyperquad} (see Lemma \\ref{lem:alg}), if $w_2(\\alpha) > r-1$, then $\\deg \\alpha = r+1$.\nIn this paper, we deal with two families for hyperquadratic continued fractions introduced in \\cite{Lasjaunias02, Lasjaunias17}.\nWe determine the exact degree of these families in particular cases.\n\nThis paper is organized as follows.\nIn Section \\ref{sec:notation}, we recall notations which are used in this paper.\nIn Section \\ref{sec:main}, we state main results of quadratic approximation for hyperquadratic continued fractions and applications of the results.\nIn Section \\ref{sec:lemma}, we give lemmas for the proof of main results.\nIn Section \\ref{sec:proof}, we prove the main results.\n\n\n\n\n\n\n\\section{Notations}\\label{sec:notation}\n\n\n\nA nonzero Laurent series $\\xi \\in \\lF_q((T ^ {- 1}))$ is represented by $\\xi = \\sum_{n = N}^{\\infty}a_n T ^ {- n}$, where $N \\in \\lZ, a_n \\in \\lF_q,$ and $a_N \\neq 0$.\nThe field $\\lF_q ((T ^ {- 1}))$ has a non-Archimedean absolute value $\\abs{\\xi} = q ^ {- N}$ and $\\abs{0} = 0$.\nThe absolute value can be uniquely extended to the algebraic closure of $\\lF_q ((T ^ {- 1}))$ and we also write $\\abs{\\cdot}$ for the extended absolute value.\n\nThe {\\itshape height} of $P(X) \\in (\\lF_q [T])[X]$, denoted by $H(P)$, is defined to be the maximum of the absolute values of the coefficients of $P(X)$.\nFor $\\alpha \\in \\overline{\\lF_q (T)}$, there exists a unique non-constant, irreducible, primitive polynomial $P(X) \\in (\\lF _q [T]) [X]$ whose leading coefficients are monic polynomials in $T$ such that $P(\\alpha) = 0$. \nThe polynomial $P(X)$ is called the {\\itshape minimal polynomial} of $\\alpha$.\nThe {\\itshape height} (resp.\\ the {\\itshape degree}, the {\\itshape inseparable degree}) of $\\alpha$, denoted by $H(\\alpha)$ (resp.\\ $\\deg \\alpha$, $\\insep \\alpha$), is defined to be the height of $P(X)$ (resp.\\ the degree of $P(X)$, the inseparable degree of $P(X)$).\nLet $n \\geq 1$ be an integer and $\\xi$ be in $\\lF_q (( T ^ {- 1}))$.\nWe denote by $w_n (\\xi)$ (resp.\\ $w_n ^ {*} (\\xi)$) the supremum of the real numbers $w$ (resp.\\ $w ^ {*}$) which satisfy \n\\begin{equation*}\n\t0 < \\abs{P(\\xi )} \\leq H(P) ^ {- w} \\quad (\\text{resp.\\ } 0 < \\abs{\\xi - \\alpha } \\leq H(\\alpha) ^ {- w ^ {*} - 1})\n\\end{equation*}\nfor infinitely many $P(X) \\in (\\lF_q [T]) [X]$ of degree at most $n$ (resp.\\ $\\alpha \\in \\overline{\\lF_q (T)}$ of degree at most $n$).\n\nLet $\\alpha \\in \\overline{\\lF_q (T)}$ be a quadratic number.\nIf $\\insep \\alpha = 1$, let $\\alpha ' \\neq \\alpha$ be the Galois conjugate of $\\alpha$.\nIf $\\insep \\alpha = 2$, let $\\alpha ' = \\alpha $.\n\nA Laurent series $\\xi \\in \\lF_q ((T ^ {- 1}))$ can be expressed as a continued fraction:\n\\begin{equation*}\n\t\\xi = a_0 + \\cfrac{1}{a_1 + \\cfrac{1}{a_2 + \\cfrac{1}{\\cdots }}},\n\\end{equation*}\nwhere $a_n \\in \\lF_q [T]$ for all $n \\geq 0$ and $\\deg _T a_m \\geq 1$ for all $m \\geq 1$.\nWe write the expansion $\\xi = [a_0, a_1, a_2, \\ldots ]$.\nThe continued fraction expansion of $\\xi$ is finite if and only if $\\xi$ is in $\\lF_q (T)$.\nIt is known that the continued fraction expansion of $\\xi$ is ultimately periodic if and only if $\\xi$ is quadratic (see Th\\'{e}or\\`{e}me 4 in \\cite[CHAPITRE IV]{Mathan70}).\nWe define sequences $(p_n)_{n \\geq -1}$ and $(q_n)_{n \\geq -1}$ by\n\\begin{equation*}\n\t\\begin{cases}\n\t\tp_{- 1} = 1, \\ p_0 = a_0, \\ p_n = a_n p_{n - 1} + p_{n - 2},\\ n \\geq 1,\\\\\n\t\tq_{- 1} = 0,\\ q_0 = 1,\\ q_n = a_n q_{n - 1} + q_{n - 2},\\ n \\geq 1.\n\t\\end{cases}\n\\end{equation*}\nWe call $(p_n \/ q_n)_{n \\geq 0}$ the {\\itshape convergent sequence} of $\\xi$ and $p_n \/ q_n$ the \\textit{$n$-th convergent} of $\\xi$.\n\nLet $n \\geq 1$ be an integer and $W, V$ be finite words.\nWe denote by $W \\bigoplus V$ the word of concatenation of $W$ and $V$.\nWe put $W ^ {[n]} := W \\bigoplus \\cdots \\bigoplus W$ ($n$ times) and $\\overline{W} := W\\bigoplus W\\bigoplus \\cdots \\bigoplus W\\bigoplus \\cdots $ (infinitely many times).\nWe write $W ^ {[0]}$ the empty word.\n\nLet $A, B$ be real numbers with $B \\neq 0$.\nWe write $A \\ll B$ (resp.\\ $A \\ll_a B$) if $\\abs{A} \\leq C \\abs{B}$ for some constant (resp.\\ some constant depending at most on $a$) $C > 0$.\nWe write $A \\asymp B$ (resp.\\ $A \\asymp_a B$) if $A \\ll B$ and $B \\ll A$ (resp.\\ $A \\ll_a B$ and $B \\ll_a A$) hold.\n\n\n\n\n\n\n\\section{Main results}\\label{sec:main}\n\n\n\nWe recall a family of hyperquadratic continued fractions.\nLet $s, t \\geq 1$ be integers and we put $q := p ^ s$ and $r := q ^ t$.\nLet $k \\geq 0$ be an integer and $\\lambda$ be in $\\lF_q ^ {*}$.\nIf $p \\neq 2$, then we assume that $\\lambda \\neq 2$ and put $\\mu := 2 - \\lambda$.\nWe define $\\Theta_k ^ {t} (\\lambda) \\in \\lF_q ((T ^ {- 1}))$ by\n\\begin{equation}\\label{eq:firstfamily}\n\t\\Theta_k ^ {t} (\\lambda)\n\t=\n\t\\begin{cases}\n\t\t[0, T ^ {[k]}, \\bigoplus_{i \\geq 1} (T, (\\lambda T, \\mu T) ^ {[(r ^ i - 1) \/ 2]}) ^ {[k+1]}] & \\text{if}\\ p \\neq 2,\\\\\n\t\t[0, T ^ {[k]} \\bigoplus_{i \\geq 1} (T, \\lambda T ^ {[r ^ i - 1]}) ^ {[k + 1]}] & \\text{if}\\ p = 2.\n\t\\end{cases}\n\\end{equation}\nLasjaunias and Ruch \\cite{Lasjaunias02} proved that $\\Theta_k ^ {t} (\\lambda)$ is hyperquadratic and satisfies the algebraic equation\n\\begin{equation}\\label{eq:firstfamilyequation}\n\t\\begin{cases}\n\t\tq_k X ^ {r + 1} - p_k X ^ r + (\\lambda \\mu ) ^ {(r - 1) \/ 2} q_{k + r} X - (\\lambda \\mu) ^ {(r - 1) \/ 2} p_{k + r} = 0 \\quad \\text{if}\\ p \\neq 2,\\\\\n\t\tq_k X ^ {r + 1} - p_k X ^ r + q_{k + r} X - p_{k + r} = 0 \\quad \\text{if}\\ p = 2,\n\t\\end{cases}\n\\end{equation}\nwhere $(p_n \/ q_n)_{n \\geq 0}$ is the convergent sequence of $\\Theta_k ^ {t} (\\lambda)$.\nNote that if $\\lambda = 1$, then $\\Theta_k ^ {t} (\\lambda)$ is ultimately periodic continued fraction, that is, $\\Theta_k ^ {t} (\\lambda)$ is quadratic.\n\nOur first main result is the following.\n\n\n\n\\begin{thm}\\label{thm:main1}\n\tLet $d \\geq 2$ be an integer and $\\Theta_k ^ {t} (\\lambda) \\in \\lF_q ((T ^ {- 1}))$ be as in \\eqref{eq:firstfamily}.\n\tAssume that $\\lambda \\neq 1$.\n\tThen we have\n\t\\begin{gather}\\label{eq:w2*}\n\t\tw_2 ^ {*} (\\Theta_k ^{t} (\\lambda)) \\geq \\max \\left\\lbrace \\frac{r - 1}{k + 1}, r - 1 - \\frac{r(r - 1)(r - 4)}{2r(k + 1) + (r - 1)(r - 2)}\\right\\rbrace ,\\\\\n\t\t\\label{eq:w2}\n\t\tw_2 (\\Theta_k ^ {t} (\\lambda)) \\geq \\max \\left\\lbrace \\frac{r - 1}{k + 1} + 1, r - \\frac{r(r - 1)(r - 3)}{2r(k + 1)+(r - 1)(r - 2)}\\right\\rbrace .\n\t\\end{gather}\n\tIf\n\t\\begin{equation}\\label{eq:maincond1}\n\t\tk \\leq \\frac{2(r - 1)}{d + \\sqrt{d ^ 2 + 4(2r - 1)d + 4}} - 1,\n\t\\end{equation}\n\tthen we have\n\t\\begin{equation}\\label{eq:mainequal1}\n\t\tw_n ^ {*} ( \\Theta_k ^ {t} (\\lambda)) = \\frac{r - 1}{k + 1}, \\quad w_n (\\Theta_k ^ {t} (\\lambda)) = \\frac{r - 1}{k + 1} + 1\n\t\\end{equation}\n\tfor all integers $2 \\leq n \\leq d$.\n\tIf\n\t\\begin{equation}\\label{eq:maincond2}\n\t\t2rd \\leq \\left( r - 2 - \\frac{r(r - 1)(r - 4)}{2r(k + 1) + (r - 1)(r - 2)}\\right) \\left( r - d - \\frac{r(r - 1)(r - 4)}{2r(k + 1) + (r - 1)(r - 2)}\\right) ,\n\t\\end{equation}\n\tthen we have\n\t\\begin{gather}\\label{eq:mainequal2}\n\t\tw_n ^ {*} (\\Theta_k ^ {t} (\\lambda)) = r - 1 - \\frac{r(r - 1)(r - 4)}{2r(k + 1) + (r - 1)(r - 2)},\\\\\n\t\t\\label{eq:mainequal22}\n\t\tw_n (\\Theta_k ^ {t} (\\lambda )) = r - \\frac{r(r - 1)(r - 3)}{2r(k + 1) + (r - 1)(r - 2)}\n\t\\end{gather}\n\tfor all integers $2 \\leq n \\leq d$.\n\\end{thm}\n\n\n\n\\begin{rem}\n\tIt is well-known and easy to see that for any $\\xi \\in \\lF_q ((T ^ {- 1}))$,\n\t\\begin{equation}\\label{eq:w1_determine}\n\t\tw_1 (\\xi) = w_1 ^ {*} (\\xi) = \\limsup_{n \\longrightarrow \\infty } \\frac{\\deg q_{n + 1}}{\\deg q_n},\n\t\\end{equation}\n\twhere $(p_n \/ q_n)_{n \\geq 0}$ is the convergent sequence of $\\xi$.\n\tSince $\\Theta_k ^ {t}(\\lambda)$ has a bounded partial quotient, we have $w_1 (\\Theta_k ^ {t}(\\lambda)) = w_1 ^ {*} (\\Theta_k ^ {t} (\\lambda)) = 1$.\n\tIf\n\t\\begin{equation*}\n\t\tr > \\frac{3d + 2 + \\sqrt{9d ^ 2 + 4d + 4}}{2},\n\t\\end{equation*}\n\tthere exists an effectively computable positive constant $C_1 (r, d)$, depending only on $r$ and $d$, such that we have \\eqref{eq:mainequal2} and \\eqref{eq:mainequal22} for all $k \\geq C_1 (r, d)$.\n\\end{rem}\n\n\n\nThe key point of the proof of Theorem \\ref{thm:main1} is to find two sequences of very good quadratic approximation, that is, ultimately periodic continued fractions.\n\nAs an application of Theorem \\ref{thm:main1}, we give a sufficient condition to determine the exact degree of $\\Theta_k ^ {t} (\\lambda)$.\n\n\n\n\\begin{cor}\\label{cor:main1}\n\tLet $\\Theta_k ^ {t} (\\lambda) \\in \\lF_q ((T ^ {- 1}))$ be as in \\eqref{eq:firstfamily}.\n\tAssume that $\\lambda \\neq 1$.\n\tIf\n\t\\begin{equation}\\label{eq:main1_cor_cond}\n\t\tk = 0 \\ \\text{or}\\ k > \\frac{(r ^ 2 - 4r + 2)(r - 1)}{2r} - 1,\n\t\\end{equation}\n\tthen we have $\\deg \\Theta_k ^ {t} (\\lambda) = r + 1$.\n\\end{cor}\n\n\n\nWe recall another family of hyperquadratic continued fractions.\nWe define a sequence $(F_n)_{n \\geq 0}$ of polynomials in $\\lF_q [T]$ by\n\\begin{equation}\\label{eq:fib}\n\tF_0 = 1, \\quad F_1 = T, \\quad F_{n + 1} = TF_n + F_{n - 1}\\quad \\text{for}\\ n \\geq 1.\n\\end{equation}\nNote that the sequence can be regarded as the analogue of the Fibonacci sequence.\nLet $s, t \\geq 1$ be integers and we put $q := 2 ^ s, r := 2 ^ t$.\nLet $\\ell \\geq 1$ be an integer, ${\\bm \\lambda} = (\\lambda_1, \\ldots , \\lambda_{\\ell}) \\in (\\lF_q ^ {*}) ^ {\\ell}, {\\bm \\varepsilon} = (\\varepsilon_1, \\varepsilon_2) \\in (\\lF_q ^ {*}) ^ 2$, and $(p_n \/ q_n)_{0 \\leq n \\leq \\ell - 1}$ be the convergent sequence of the continued fraction $[\\lambda_1 T, \\ldots , \\lambda_{\\ell}T]$.\nWe consider the algebraic equation\n\\begin{equation}\\label{eq:char2eq}\n\tq_{\\ell - 1} X ^ {r + 1} + p_{\\ell - 1} X ^ r + (\\varepsilon_1 q_{\\ell - 2} F_{r-1} + \\varepsilon_2 q_{\\ell - 1} F_{r - 2}) X + \\varepsilon_1 p_{\\ell - 2} F_{r - 1} + \\varepsilon_2 p_{\\ell - 1} F_{r - 2} = 0.\n\\end{equation}\nLasjaunias \\cite{Lasjaunias17} proved that \\eqref{eq:char2eq} has a unique root $\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})$ in $\\lF_q ((T ^ {- 1}))$ with $\\abs{\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})} \\geq q$, $\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})$ is hyperquadratic and its continued fraction expansion is the following:\n\\begin{equation}\\label{eq:secondfamily}\n\t\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon}) = [\\lambda_1 T, \\lambda_2 T, \\ldots ],\n\\end{equation}\nwhere a sequence $(\\lambda_n)_{n \\geq 1}$ satisfy\n\\begin{gather*}\n\t\\lambda_{\\ell + r m + 1} = (\\varepsilon_2 \/ \\varepsilon_1) \\varepsilon_2 ^ {(- 1) ^ {m + 1}} \\lambda_{m + 1} ^ r,\\\\\n\t\\lambda_{\\ell + r m + i} = (\\varepsilon_1 \/ \\varepsilon_2 ) ^ {(- 1) ^ i} \\quad \\text{for}\\ m \\geq 0 \\ \\text{and}\\ 2 \\leq i \\leq r.\n\\end{gather*}\nNote that if $\\varepsilon_1 = \\varepsilon_2 = 1$ and $\\lambda_j = 1$ for all $1\\leq j \\leq \\ell$, then $\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})$ is ultimately periodic continued fraction, that is, $\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})$ is quadratic.\n\nIn this paper, we consider the continued fractions $\\Phi_{\\ell } ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})$ in particular cases.\n\n\n\\begin{thm}\\label{thm:main2}\n\tLet $\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})$ be as in \\eqref{eq:secondfamily}, $m \\geq 1, d \\geq 2$ be integers with $m \\leq \\ell$, and $\\lambda$ be in $\\lF_q ^ {*}$ with $\\lambda \\neq 1$.\n\t\\begin{enumerate}\n\t\t\\item[(1)] Assume that $\\lambda _1 = \\lambda , \\lambda _i = 1$ for all $2 \\leq i \\leq m, \\varepsilon_1 = \\varepsilon_2 = 1$ and if $m < \\ell $, then $\\lambda_{m + 1} \\neq 1$.\n\t\tThen we have\n\t\t\\begin{equation}\\label{eq:main22}\n\t\t\tw_2 ^ {*} (\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon}))\\geq \\frac{m}{\\ell}(r - 1),\\quad \n\t\t\tw_2 (\\Phi_{\\ell} ^ t ({\\bm \\lambda }; {\\bm \\varepsilon})) \\geq 1 + \\frac{m}{\\ell} (r - 1).\n\t\t\\end{equation}\n\t\tIf\n\t\t\\begin{equation}\\label{eq:maincond3}\n\t\t\t\\frac{\\ell}{m} \\leq \\frac{2(r - 1)}{d + \\sqrt{d ^ 2 + 4(2r - 1)d + 4}},\n\t\t\\end{equation}\n\t\tthen we have\n\t\t\\begin{equation}\\label{eq:mainequal3}\n\t\t\tw_n ^ {*} (\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})) = \\frac{m}{\\ell}(r - 1), \\quad \n\t\t\tw_n (\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})) = 1 + \\frac{m}{\\ell}(r - 1)\n\t\t\\end{equation}\n\t\tfor all integers $2 \\leq n \\leq d$.\n\t\t\n\t\t\\item[(2)] Assume that $r$ is a power of $q, \\varepsilon_1 = \\varepsilon_2 = 1$ and $\\lambda_i = \\lambda$ for all $1 \\leq i \\leq \\ell$.\n\t\tThen we have\n\t\t\\begin{gather}\\label{eq:mainequal6}\n\t\t\tw_2 ^ {*} (\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})) \\geq r - 1 - \\frac{r(r - 1)(r - 4)}{2\\ell r + (r - 1)(r - 2)},\\\\\n\t\t\t\\label{eq:mainequal66}\n\t\t\tw_2 (\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})) \\geq r - \\frac{r(r - 1)(r - 3)}{2\\ell r + (r - 1)(r - 2)}.\n\t\t\\end{gather}\n\t\tIf\n\t\t\\begin{equation}\\label{eq:maincond7}\n\t\t\t2 r d \\leq \\left( r - 2 - \\frac{r(r - 1)(r - 4)}{2 \\ell r + (r - 1)(r - 2)}\\right) \\left( r - d - \\frac{r(r - 1)(r - 4)}{2 \\ell r + (r - 1)(r - 2)}\\right) ,\n\t\t\\end{equation}\n\t\tthen we have\n\t\t\\begin{gather}\\label{eq:mainequal7}\n\t\t\tw_n ^ {*} (\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})) = r - 1 - \\frac{r(r - 1)(r - 4)}{2 \\ell r + (r - 1)(r - 2)},\\\\\n\t\t\t\\label{eq:mainequal77}\n\t\t\tw_n (\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})) = r - \\frac{r (r - 1)(r - 3)}{2 \\ell r + (r - 1)(r - 2)}\n\t\t\\end{gather}\n\t\tfor all integers $2 \\leq n \\leq d$.\n\t\\end{enumerate}\n\\end{thm}\n\n\n\\begin{rem}\n\tBy \\eqref{eq:w1_determine}, we have $w_1 (\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})) = w_1 ^ {*} (\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})) = 1$.\n\tIf\n\t\\begin{equation*}\n\t\tr > \\frac{3d + 2 + \\sqrt{9d ^ 2 + 4d + 4}}{2},\n\t\\end{equation*}\n\tthere exists an effectively computable positive constant $C_2 (r, d)$, depending only on $r$ and $d$, such that we have \\eqref{eq:mainequal7} and \\eqref{eq:mainequal77} for all $\\ell \\geq C_2 (r, d)$.\n\\end{rem}\n\n\n\n\\begin{cor}\\label{cor:main2}\n\tLet $\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})$ be as in \\eqref{eq:secondfamily}, $m \\geq 1, d \\geq 2$ be integers with $m \\leq \\ell$, and $\\lambda$ be in $\\lF_q ^ {*}$ with $\\lambda \\neq 1$.\n\t\\begin{enumerate}\n\t\t\\item[(1)] Assume that $\\lambda _1 = \\lambda , \\lambda _i = 1$ for all $2 \\leq i \\leq m, \\varepsilon_1 = \\varepsilon_2 = 1$ and if $m < \\ell $, then $\\lambda_{m + 1} \\neq 1$.\n\t\tIf\n\t\t\\begin{equation*}\n\t\t\t\\frac{m}{\\ell} > \\frac{r - 2}{r - 1},\n\t\t\\end{equation*}\n\t\tthen we have $\\deg \\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon}) = r + 1$.\n\t\t\n\t\t\\item[(2)] Assume that $r$ is a power of $q, \\varepsilon_1 = \\varepsilon_2 = 1$ and $\\lambda_i = \\lambda$ for all $1 \\leq i \\leq \\ell$.\n\t\tIf\n\t\t\\begin{equation*}\n\t\t\t\\ell > \\frac{(r ^ 2 - 4r + 2)(r - 1)}{2r},\n\t\t\\end{equation*}\n\t\tthen we have $\\deg \\Phi_{\\ell} ^ t({\\bm \\lambda}; {\\bm \\varepsilon}) = r + 1$.\n\t\\end{enumerate}\n\\end{cor}\n\n\n\nIn the last part of this section, we mention a problem associated to Problem \\ref{prob:main1}, Corollary \\ref{cor:main1} and \\ref{cor:main2}.\n\n\n\n\\begin{prob}\n\tLet $n\\geq 2$ be an integer.\n\tDetermine the set of all values taken by $w_n - w_n ^ {*}$ over the set of algebraic Laurent series.\n\\end{prob}\n\n\n\n\n\n\n\\section{Preliminaries}\\label{sec:lemma}\n\n\n\nIn this section, we gather lemmas for the proof of main results.\n\n\n\n\\begin{lem}\\label{lem:alg}{\\rm (\\cite[Theorem 5.2]{Ooto17}).}\n\tLet $n \\geq 1$ be an integer and $\\alpha \\in \\lF_q ((T ^ {- 1}))$ be an algebraic Laurent series.\n\tThen we have \n\t\\begin{equation*}\n\t\tw_n ^ {*} (\\alpha), w_n (\\alpha) \\leq \\deg \\alpha - 1.\n\t\\end{equation*}\n\\end{lem}\n\n\n\nLemma \\ref{lem:contidif} and \\ref{lem:conti.lem} are immediately seen.\n\n\n\n\\begin{lem}\\label{lem:contidif}\n\tLet $\\xi = [0, a_1, a_2, \\ldots ], \\zeta =[0, b_1, b_2, \\ldots ]$ be in $\\lF_q ((T ^ {- 1}))$.\n\tAssume that there exists an integer $k \\geq 1$ such that $a_n = b_n$ for all $1 \\leq n \\leq k$ and $a_{k + 1} \\neq b_{k + 1}$.\n\tThen we have\n\t\\begin{equation*}\n\t|\\xi - \\zeta| = \\frac{\\abs{a_{k + 1} - b_{k + 1}}}{\\abs{a_{k + 1} b_{k + 1}}\\abs{q_k} ^ 2},\n\t\\end{equation*}\n\twhere $(p_n \/ q_n)_{n \\geq 0}$ is the convergent sequence of $\\xi$.\n\\end{lem}\n\n\n\n\\begin{lem}\\label{lem:conti.lem}\n\tLet $\\xi = [a_0, a_1, \\ldots ]$ be in $\\lF_q ((T ^ {- 1}))$ and $(p_n \/ q_n)_{n \\geq 0}$ be the convergent sequence of $\\xi$.\n\tThen, for any $n \\geq 1$, we have\n\t\\begin{equation}\\label{enum:q}\n\t\t|q_n| = |a_1| |a_2| \\cdots |a_n|.\n\t\\end{equation}\n\\end{lem}\n\n\n\n\\begin{lem}\\label{lem:conj2}{\\rm (\\cite[Lemma 4.6]{Ooto17}).}\n\tLet $r, s \\geq 1$ be integers and $\\alpha = [0, a_1,\\ldots ,a_r, \\overline{a_{r + 1}, \\ldots , a_{r + s}}] \\in \\lF_q ((T ^ {- 1}))$ be an ultimately periodic continued fraction with $a_r \\neq a_{r + s}$.\n\tLet $(p_n \/ q_n)_{n \\geq 0}$ be the convergent sequence of $\\alpha$.\n\tThen we have\n\t\\begin{equation*}\n\t\t\\frac{\\min (|a_r|, |a_{r + s}| )}{|q_r| ^ 2} \\leq |\\alpha - \\alpha '| \\leq \\frac{|a_r a_{r + s}|}{|q_r| ^ 2}.\n\t\\end{equation*}\n\\end{lem}\n\n\n\nThe following lemma is well-known and easy to see.\n\n\n\n\\begin{lem}\\label{lem:quadheightupper}\n\tLet $r \\geq 0, s \\geq 1$ be integers and $\\alpha = [0, a_1, \\ldots , a_r, \\overline{a_{r + 1}, \\ldots , a_{r + s}}] \\in \\lF_q ((T ^ {- 1}))$ be an ultimately periodic continued fraction.\n\tLet $(p_n \/ q_n)_{n \\geq 0}$ be the convergent sequence of $\\xi$.\n\tThen we have $H(\\alpha) \\leq |q_r q_{r + s}|$.\n\\end{lem}\n\n\n\n\\begin{lem}\\label{lem:quadheight}\n\tLet $b, c, d \\in \\lF_q [T]$ be distinct polynomials with $\\deg_T b, \\deg_T c, \\deg_T d \\geq 1$.\n\tLet $n, m, \\ell \\geq 1$ be integers and $a_1, \\ldots , a_{n - 1} \\in \\lF_q [T]$ be polynomials with $\\deg_T a_i \\geq 1$ for all $1 \\leq i \\leq n - 1$.\n\tWe put\n\t\\begin{gather*}\n\t\t\\alpha_1 := [0, a_1, \\ldots , a_{n - 1}, c, \\overline{b}], \\quad \n\t\t\\alpha_2 := [0, a_1, \\ldots , a_{n - 1}, c, \\overline{b ^ {[m]}, d}],\\\\\n\t\t\\alpha_3 := [0, a_1, \\ldots , a_{n - 1}, c, \\overline{b ^ {[m]}, c, b ^ {[\\ell]}}], \\quad\n\t\t\\alpha_4 := [0, a_1, \\ldots , a_{n - 1}, c, \\overline{(b, d) ^ {[m]}, c, (b, d) ^ {[\\ell]}}].\n\t\\end{gather*}\n\tLet $(p_{i, k} \/ q_{i, k})_{k \\geq 0}$ be the convergent sequence of $\\xi_i$ for $i = 1, 2, 3, 4$.\n\tThen we have\n\t\\begin{gather}\\label{eq:quadheight12}\n\t\tH(\\alpha_1) \\asymp_{b, c} \\abs{q_{1, n}} ^ 2, \\quad\n\t\tH(\\alpha_2) \\asymp _{b, c, d} \\abs{q_{2, n} q_{2, n + m + 1}},\\\\\n\t\t\\label{eq:quadheight3}\n\t\tH(\\alpha_3) \\asymp_{b, c} \\abs{q_{3, n} q_{3, n + m + \\ell + 1}}, \\quad\n\t\tH(\\alpha_4) \\asymp_{b, c, d} \\abs{q_{4, n} q_{4, n + 2m + 2\\ell + 1}}.\n\t\\end{gather}\n\\end{lem}\n\n\n\n\\begin{proof}\n\tSee Lemma 4.7 in \\cite{Ooto17} for the proof of \\eqref{eq:quadheight12}.\n\t\n\tIt follows from Lemma \\ref{lem:quadheightupper} that $H(\\alpha_3) \\leq |q_{3, n} q_{3, n + m + \\ell + 1}|$.\n\tWe put \n\t\\begin{equation*}\n\t\t\\beta_3 :=[0, a_1, \\ldots , a_{n - 1}, c, b ^ {[m]}, c, \\overline{b}].\n\t\\end{equation*}\n\t\n\tLet $P_3(X)$ and $Q_3(X)$ be the minimal polynomial of $\\alpha_3$ and $\\beta_3$, respectively.\n\tSince $P_3$ and $Q_3$ do not have a common root, we have\n\t\\begin{equation*}\n\t\t1 \\leq |\\Res(P_3, Q_3)| \\leq H(\\beta_3) ^ 2 H(\\alpha_3) ^ 2 |\\alpha_3 - \\beta_3| |\\alpha_3' - \\beta_3| |\\alpha_3 - \\beta_3'| |\\alpha_3' - \\beta_3'|.\n\t\\end{equation*}\n\tBy Lemma \\ref{lem:contidif}, \\ref{lem:conti.lem} and \\ref{lem:conj2}, we obtain\n\t\\begin{gather*}\n\t\t|\\alpha_3 - \\beta_3| \\ll_{b, c} |q_{3, n + 2m + \\ell + 1}| ^ {- 2}, \\quad \n\t\t|\\alpha_3 - \\beta_3'| \\ll_{b, c} |q_{3, n + m + 1}| ^ {- 2},\\\\\n\t\t|\\alpha_3' - \\beta_3|, |\\alpha_3' - \\beta_3'| \\ll_{b, c} |q_{3, n}| ^ {- 2}.\n\t\\end{gather*}\n\tTherefore, by Lemma \\ref{lem:conti.lem} and \\ref{lem:quadheightupper}, we have $H(\\alpha_3) \\gg_{b, c} |q_{3, n} q_{3, n + m + \\ell + 1}|$.\n\t\n\tIn the same way to the above proof, we obtain $H(\\alpha_4) \\asymp_{b, c, d} \\abs{q_{4, n} q_{4, n + 2m + 2\\ell + 1}}$.\n\\end{proof}\n\n\n\n\\begin{lem}\\label{lem:wnwn*}{\\rm (\\cite[Proposition 5.6]{Ooto17}).}\n\tLet $n \\geq 1$ be an integer and $\\xi \\in \\lF_q ((T ^ {- 1}))$.\n\tThen we have \n\t\\begin{equation*}\n\t\tw_n ^ {*} (\\xi) \\leq w_n (\\xi).\n\t\\end{equation*}\n\\end{lem}\n\n\n\n\\begin{lem}\\label{lem:wnlower}\n\tLet $n, m \\geq 1$ be integers with $n \\geq m$ and $\\xi \\in \\lF_q ((T ^ {- 1}))$ be not algebraic of degree at most $m$.\n\tThen we have $w_n (\\xi) \\geq m$.\n\\end{lem}\n\n\n\n\\begin{proof}\n\tSince the sequence $(w_k (\\xi))_{k \\geq 1}$ is increasing, we have $w_n (\\xi) \\geq m$ (see e.g.\\ \\cite[page~200]{Bugeaud04}).\n\\end{proof}\n\n\n\nThe following lemma is well-known and immediately seen.\n\n\n\n\\begin{lem}\\label{lem:height}\n\tLet $P(X)$ be in $(\\lF_q [T]) [X]$.\n\tAssume that $P(X)$ can be factorized as\n\t\\begin{equation*}\n\t\tP(X) = A \\prod_{i = 1}^{n} (X - \\alpha_i),\n\t\\end{equation*}\n\twhere $A \\in \\lF_q [T]$ and $\\alpha_i \\in \\overline{\\lF_q (T)}$ for $1 \\leq i \\leq n$.\n\tThen we have\n\t\\begin{equation*}\n\t\tH(P) = |A| \\prod_{i = 1}^{n} \\max (1, |\\alpha_i|).\n\t\\end{equation*}\n\\end{lem}\n\n\n\nThe following lemma is immediately seen by the definition of discriminant.\n\n\n\n\\begin{lem}\\label{lem:Galois}\n\tLet $\\alpha \\in \\overline{\\lF_q (T)}$ be a quadratic number.\n\tIf $\\alpha \\neq \\alpha'$, then we have\n\t\\begin{equation*}\n\t\t|\\alpha - \\alpha'| \\geq H(\\alpha) ^ {- 1}.\n\t\\end{equation*}\n\\end{lem}\n\n\n\nThe following two lemmas are key lemmas for proving Theorem \\ref{thm:main1} and \\ref{thm:main2}.\n\n\n\n\\begin{lem}\\label{lem:bestquad}{\\rm (\\cite[Lemma 3.19]{Ooto18}).}\n\tLet $d \\geq 2$ be an integer.\n\tLet $\\xi$ be in $\\lF_q ((T ^ {- 1}))$, $\\theta , \\delta$ be positive numbers, and $\\varepsilon$ be a non-negative number.\n\tAssume that there exist a sequence $(\\alpha_j)_{j \\geq 1}$ and positive number $c$ such that for any $j \\geq 1$, $\\alpha_j \\in \\overline{\\lF_q (T)}$ is quadratic with $0 < |\\alpha_j - \\alpha_j'| \\leq c$ and $\\xi \\neq \\alpha_j$, $(H(\\alpha_j))_{j \\geq 1}$ is a divergent increasing sequence, and\n\t\\begin{gather*}\n\t\t\\limsup_{k \\rightarrow \\infty } \\frac{\\log H(\\alpha_{k + 1})}{\\log H(\\alpha_k)} \\leq \\theta ,\\\\\n\t\t\\lim_{k \\rightarrow \\infty } \\frac{- \\log |\\xi - \\alpha_k|}{\\log H(\\alpha_k)} = d + \\delta , \\quad \\lim_{k \\rightarrow \\infty } \\frac{- \\log |\\alpha_k - \\alpha_k'|}{\\log H(\\alpha_k)} = \\varepsilon .\n\t\\end{gather*}\n\tIf $2 d \\theta \\leq (d - 2 + \\delta ) \\delta $, then we have for all $2 \\leq n \\leq d$,\n\t\\begin{equation}\\label{eq:bestquad_eq}\n\t\tw_n ^ {*} (\\xi) = d - 1 + \\delta , \\quad w_n(\\xi) = d - 1 + \\delta + \\varepsilon .\n\t\\end{equation}\n\\end{lem}\n\n\n\n\\begin{lem}\\label{lem:bestquad2}\n\tLet $\\xi$ be in $\\lF_q ((T ^ {- 1})) \\setminus \\lF_q (T)$.\n\tAssume that there exists a sequence $(\\alpha_j)_{j \\geq 1}$ such that for any $j \\geq 1$, $\\alpha_j \\in \\overline{\\lF_q (T)}$ is quadratic with $\\alpha_j \\neq \\alpha_j'$ and $\\xi \\neq \\alpha_j$, and $(H(\\alpha_j))_{j \\geq 1}$ is a divergent increasing sequence.\n\tIf there exist limits of the sequences \n\t\\begin{equation*}\n\t\t\\left( \\frac{- \\log |\\xi - \\alpha_j|}{\\log H(\\alpha_j)}\\right) _{j \\geq 1}\\ \\text{and}\\ \\left( \\frac{- \\log |\\alpha_j - \\alpha_j'|}{\\log H(\\alpha_j)}\\right) _{j \\geq 1},\n\t\\end{equation*}\n\tthen we have\n\t\\begin{gather}\\label{eq:quad1}\n\t\tw_2 ^ {*} (\\xi) \\geq \\lim_{j \\rightarrow \\infty } \\frac{- \\log |\\xi - \\alpha_j|}{\\log H(\\alpha_j)} - 1,\\\\\n\t\t\\label{eq:quad2}\n\t\tw_2 (\\xi) \\geq \\lim_{j \\rightarrow \\infty } \\frac{- \\log |\\xi - \\alpha_j|}{\\log H(\\alpha_j)} + \\lim_{j \\rightarrow \\infty }\\frac{- \\log |\\alpha_j - \\alpha_j'|}{\\log H(\\alpha_j)} - 1.\n\t\\end{gather}\n\\end{lem}\n\n\n\n\\begin{proof}\n\tThe inequality \\eqref{eq:quad1} is immediately seen.\n\tIn what follows, we show \\eqref{eq:quad2}.\n\tWe put\n\t\\begin{equation*}\n\t\t\\gamma_1 := \\lim_{j \\rightarrow \\infty } \\frac{- \\log |\\xi - \\alpha_j|}{\\log H(\\alpha_j)}, \\quad \\gamma_2 := \\lim_{j \\rightarrow \\infty } \\frac{- \\log |\\alpha_j - \\alpha_j'|}{\\log H(\\alpha_j)}.\n\t\\end{equation*}\n\t\n\tWe may assume that $\\gamma_1 > 1 $ and $\\gamma_2 > 0$.\n\tIn fact, if $\\gamma_2 \\leq 0$, then we have \\eqref{eq:quad2} by Lemma \\ref{lem:wnwn*}.\n\tBy Lemma \\ref{lem:Galois}, we have $\\gamma_2 \\leq 1$, which implies $\\gamma_2 < \\gamma_1$.\n\tIf $\\gamma_1 \\leq 1$, then we obtain \\eqref{eq:quad2} by Lemma \\ref{lem:wnlower}.\n\t\n\tThen we obtain\n\t\\begin{equation*}\n\t\t|\\xi - \\alpha_j| < |\\alpha_j - \\alpha_j'| \\leq 1\n\t\\end{equation*}\n\tfor all sufficiently large $j$.\n\tTherefore, we have\n\t\\begin{gather*}\n\t\t\\max (1, |\\xi|) = \\max (1, |\\alpha_j|) = \\max (1, |\\alpha_j'|),\\\\\n\t\t|\\xi - \\alpha_j'| = |\\alpha_j - \\alpha_j'|\n\t\\end{gather*}\n\tfor all sufficiently large $j$.\n\tIt follows from Lemma \\ref{lem:height} that\n\t\\begin{equation*}\n\t\tH(P_j) = |A_j| \\max (1, |\\xi|) ^ 2\n\t\\end{equation*}\n\tfor all sufficiently large $j$, where $P_j (X) = A_j (X - \\alpha_j) (X - \\alpha_j ')$ is the minimal polynomial of $\\alpha_j$.\n\tHence, we have\n\t\\begin{equation*}\n\t\t|P_j (\\xi)| = H(P_j) |\\xi - \\alpha_j| |\\alpha_j - \\alpha_j'| \\max (1, |\\xi|) ^ {- 2}\n\t\\end{equation*}\n\tfor all sufficiently large $j$.\n\tThus, we obtain\n\t\\begin{equation*}\n\t\t\\lim_{j\\rightarrow \\infty }\\frac{- \\log |P_j (\\xi)|}{\\log H(P_j)} = \\gamma_1 + \\gamma_2 - 1,\n\t\\end{equation*}\n\twhich implies \\eqref{eq:quad2}.\n\\end{proof}\n\n\n\n\n\n\n\\section{Proof of Main results}\\label{sec:proof}\n\n\n\nIn this section, we prove the main results, that is, Theorem \\ref{thm:neq_alg}, \\ref{thm:main1}, \\ref{thm:main2}, Corollary \\ref{cor:main1} and \\ref{cor:main2}.\n\n\n\n\\begin{proof}[Proof of Theorem \\ref{thm:main1}]\n\tFor an integer $n \\geq 3$, we put\n\t\\begin{gather*}\n\t\t\\alpha_n :=\n\t\t\\begin{cases}\n\t\t\t[0, T ^ {[k]}, \\bigoplus_{1 \\leq i \\leq n - 1}(T, (\\lambda T, \\mu T) ^ {[(r ^ i - 1) \/ 2]}) ^ {[k + 1]}, T, \\overline{\\lambda T, \\mu T}] & \\text{if}\\ p \\neq 2,\\\\\n\t\t\t[0, T ^ {[k]}, \\bigoplus_{1 \\leq i \\leq n - 1}(T, \\lambda T ^ {[r ^ i - 1]}) ^ {[k + 1]}, T, \\overline{\\lambda T}] & \\text{if}\\ p = 2,\n\t\t\\end{cases}\\\\\n\t\t\\beta_n :=\n\t\t\\begin{cases}\n\t\t\t[0, T ^ {[k]}, \\bigoplus_{1 \\leq i \\leq n - 1} (T, (\\lambda T, \\mu T) ^ {[(r ^ i - 1) \/ 2]}) ^ {[k + 1]}, \\overline{T, (\\lambda T, \\mu T) ^ {[(r ^ n - 1) \/ 2]}}] & \\text{if}\\ p \\neq 2,\\\\\n\t\t\t[0, T ^ {[k]}, \\bigoplus_{1 \\leq i \\leq n - 1} (T, \\lambda T ^ {[r ^ i - 1]}) ^ {[k + 1]}, \\overline{T, \\lambda T ^ {[r ^ n - 1]}}] & \\text{if}\\ p = 2.\n\t\t\\end{cases}\n\t\\end{gather*}\n\tThen we have\n\t\\begin{align*}\n\t\t\\beta_n =\n\t\t[0, T ^ {[k]}, \\bigoplus_{1 \\leq i \\leq n - 2}(T, (\\lambda T, & \\mu T) ^ {[(r ^ i - 1) \/ 2]}) ^ {[k + 1]}, (T, (\\lambda T, \\mu T) ^ {[(r ^ {n - 1} - 1) \/ 2]})^{[k]},\\\\\n\t\t& T, \\overline{(\\lambda T, \\mu T) ^ {[(r ^ {n - 1} - 1) \/ 2]}, T, (\\lambda T, \\mu T) ^ {[(r ^ n - r ^ {n - 1}) \/ 2]}}]\t\n\t\\end{align*}\n\tif $p \\neq 2$, and\n\t\\begin{align*}\n\t\t\\beta_n\t=\n\t\t[0, T ^ {[k]}, \\bigoplus_{1 \\leq i \\leq n - 2} (T, \\lambda T ^ {[r ^ i - 1]}) ^ {[k + 1]}, (T, \\lambda T ^ {[r ^ {n - 1} - 1]}) ^ {[k]}, T, \\overline{\\lambda T ^ {[r ^ {n - 1} - 1]}, T, \\lambda T ^ {[r ^ n - r ^ {n - 1}]}}]\n\t\\end{align*}\n\tif $p = 2$.\n\tIt follows from Lemma \\ref{lem:conti.lem} and \\ref{lem:quadheight} that\n\t\\begin{equation*}\n\t\tH(\\alpha_n) \\asymp q ^ {2 (k + 1) (r ^ n - 1) \/ (r - 1)}, \\quad\n\t\tH(\\beta_n) \\asymp q ^ {2 (k + 1) (r ^ n - 1) \/ (r - 1) + r ^ n - 2 r ^ {n - 1}}.\n\t\\end{equation*}\n\tSince $\\Theta_k ^ {t} (\\lambda)$ and $\\alpha_n$ have the same first $((k + 1) \\sum_{i = 0} ^ {n - 1} r ^ i + r ^ n - 1)$-th partial quotients, while the next partial quotient are different, we have\n\t\\begin{equation*}\n\t\t\\abs{\\Theta_k ^ {t} (\\lambda) - \\alpha_n} \\asymp q ^ {- 2 (k + 1) (r ^ n - 1) \/ (r - 1) - 2 r ^ n}\n\t\\end{equation*}\n\tby Lemma \\ref{lem:contidif} and \\ref{lem:conti.lem}.\n\tSimilary, since $\\Theta_k ^ {t} (\\lambda)$ and $\\beta_n$ have the same first $((k + 1) \\sum_{i = 0} ^ {n} r ^ i + r ^ n - 1)$-th partial quotients, while the next partial quotient are different, we get\n\t\\begin{equation*}\n\t\t\\abs{\\Theta_k ^ {t} (\\lambda) - \\beta_n} \\asymp q ^ {- 2 (k + 1) (r ^ {n + 1} - 1) \/ (r - 1) - 2 r ^ n}.\n\t\\end{equation*}\n\tBy Lemma \\ref{lem:conti.lem} and \\ref{lem:conj2}, we obtain\n\t\\begin{equation*}\n\t\t\\abs{\\alpha_n - \\alpha_n'} \\asymp q ^ {- 2 (k + 1) (r ^ n - 1) \/ (r - 1)},\\quad\n\t\t\\abs{\\beta_n - \\beta_n'} \\asymp q ^ {- 2 (k + 1) (r ^ n - 1) \/ (r - 1) + 2 r ^ {n - 1}}.\n\t\\end{equation*}\n\tTherefore, we deduce that\n\t\\begin{gather*}\n\t\t\\lim_{n \\rightarrow \\infty } \\frac{\\log H(\\alpha_{n + 1})}{\\log H(\\alpha_n)} = \\lim_{n \\rightarrow \\infty } \\frac{\\log H(\\beta_{n + 1})}{\\log H(\\beta_n)} = r,\\\\\n\t\t\\lim_{n \\rightarrow \\infty } \\frac{- \\log \\abs{\\Theta_k ^ {t} (\\lambda) - \\alpha_n}}{\\log H(\\alpha_n)} = 1 + \\frac{r - 1}{k + 1},\\quad \n\t\t\\lim_{n \\rightarrow \\infty } \\frac{- \\log \\abs{\\alpha_n - \\alpha_n'}}{\\log H(\\alpha_n)} = 1,\\\\\n\t\t\\lim_{n \\rightarrow \\infty } \\frac{- \\log \\abs{\\Theta_k ^ {t} (\\lambda) - \\beta_n}}{\\log H(\\beta_n)} = r - \\frac{r (r - 1)(r - 4)}{2 r (k + 1) + (r - 1) (r - 2)},\\\\\n\t\t\\lim_{n \\rightarrow \\infty } \\frac{- \\log \\abs{\\beta_n - \\beta_n'}}{\\log H(\\beta_n)} = 1 - \\frac{r (r - 1)}{2 r (k + 1) + (r - 1)(r - 2)}.\n\t\\end{gather*}\n\tHence, we obtain \\eqref{eq:w2*} and \\eqref{eq:w2} by Lemma \\ref{lem:bestquad2}.\n\t\n\tIf \\eqref{eq:maincond1} holds, then we heve\n\t\\begin{equation*}\n\t\t2 r d \\leq \\left( \\frac{r - 1}{k + 1} - 1\\right) \\left( \\frac{r - 1}{k + 1} - d + 1 \\right) .\n\t\\end{equation*}\n\tTherefore, by Lemma \\ref{lem:bestquad}, we obtain \\eqref{eq:mainequal1}.\n\t\n\tSimilaly, if \\eqref{eq:maincond2} holds, then we heve \\eqref{eq:mainequal2} and \\eqref{eq:mainequal22} by Lemma \\ref{lem:bestquad}.\n\\end{proof}\n\n\n\n\\begin{proof}[Proof of Corollary \\ref{cor:main1}]\n\tIt follows from \\eqref{eq:hyperquad} that $\\deg \\Theta_k ^ {t} (\\lambda) \\leq r+1$.\n\tIf \\eqref{eq:main1_cor_cond} holds, then we have $w_2(\\Theta_k ^ {t} (\\lambda)) > r-1$ by Theorem \\ref{thm:main1}.\n\tTherefore, by Lemma \\ref{lem:alg}, we obtain $\\deg \\Theta_k ^ {t} (\\lambda) = r+1$.\n\\end{proof}\n\n\n\n\\begin{proof}[Proof of Theorem \\ref{thm:neq_alg}]\n\tIt follows from $q \\geq 4$ that we can take $\\lambda \\in \\lF_q ^ {*}$ with $\\lambda \\neq 1$ and $\\lambda \\neq 2$.\n\tLet $n \\geq 2$ be integers.\n\tSince $r \\geq (3n + 2 + \\sqrt{9n^2 + 4n + 4})\/2$, we take $r = p^t$, where $t \\geq 0$ is an integer with\n\t\\begin{equation*}\n\t\t0 \\leq \\frac{2(r - 1)}{n + \\sqrt{n ^ 2 + 4(2r - 1)n + 4}} - 1.\n\t\\end{equation*}\n\tThen, by Theorem \\ref{thm:main1} and Corollary \\ref{cor:main1}, we have $w_n(\\Theta_0 ^ {t} (\\lambda)) \\neq w_n^{*}(\\Theta_0 ^ {t} (\\lambda))$ and $\\deg \\Theta_0 ^ {t} (\\lambda) = r + 1$.\n\\end{proof}\n\n\n\n\\begin{proof}[Proof of Theorem \\ref{thm:main2}]\n\tSince $w_n (\\xi) = w_n (\\xi^{- 1})$ and $w_n ^ {*} (\\xi) = w_n ^ {*} (\\xi ^{- 1})$ for all $0 \\neq \\xi \\in \\lF_q ((T ^ {- 1}))$ and $n \\geq 1$, we consider $\\Phi := \\Phi_{\\ell} ^ t({\\bm \\lambda}; {\\bm \\varepsilon}) ^ {- 1} = [0, \\lambda_1T, \\lambda_2T, \\ldots ]$ instead of $\\Phi_{\\ell} ^ t ({\\bm \\lambda}; {\\bm \\varepsilon})$.\n\tBy the assumption, we obtain\n\t\\begin{equation*}\n\t\t\\lambda_n\n\t\t=\n\t\t\\begin{cases}\n\t\t\t\\lambda_{i + 1} ^ {r ^ j} & \\text{if}\\ n = 1 + \\ell \\sum_{h = 0} ^ {j - 1} r ^ h + i r ^ j \\ \\text{for some}\\ j \\geq 0, 0 \\leq i < \\ell ,\\\\\n\t\t\t1 & \\text{otherwise}.\n\t\t\\end{cases}\n\t\\end{equation*}\n\t\n\tFirst, we prove (1).\n\tFor an integer $n \\geq 1$, we put\n\t\\begin{equation*}\n\t\t\\alpha_n := [0, \\lambda_1T, \\lambda_2T, \\ldots ,\\lambda_{i(n)}T, \\overline{T}],\n\t\\end{equation*}\n\twhere $i(n) = 1 + \\ell \\sum_{j = 0} ^ {n - 1} r ^ j$.\n\tThen we have $\\lambda_{i(n)} = \\lambda ^ {r ^ n}$.\n\tIn a similar way to the proof of Theorem \\ref{thm:main1}, we obtain\n\t\\begin{gather*}\n\t\\abs{\\Phi - \\alpha_n} \\asymp q ^ {- 2 \\ell (r ^ n - 1) \/ (r - 1) - 2 m r ^ n},\\quad \n\t\\abs{\\alpha_n - \\alpha_n'} \\asymp q ^ {- 2 \\ell (r ^ n - 1) \/ (r - 1)},\\\\\n\tH(\\alpha_n) \\asymp q ^ {2 \\ell (r ^ n - 1) \/ (r - 1)}.\n\t\\end{gather*}\n\tTherefore, we have \\eqref{eq:main22} by Lemma \\ref{lem:bestquad2}.\n\tIf \\eqref{eq:maincond3} holds, then we obtain \\eqref{eq:mainequal3} by Lemma \\ref{lem:bestquad}.\n\t\n\tNext, we prove (2).\n\tFor an integer $n \\geq 1$, we put\n\t\\begin{equation*}\n\t\\beta_n := [0, \\lambda_1T, \\lambda_2T, \\ldots ,\\lambda_{i(n) - 1}T, \\overline{\\lambda ^ {r ^ n}T, T ^ {[r ^ n - 1]}}],\n\t\\end{equation*}\n\tThen we have\n\t\\begin{equation*}\n\t\\beta_n = [0, \\lambda_1T, \\lambda_2T, \\ldots ,\\lambda_{i(n) - r ^ {n - 1}}T, \\overline{T ^ {[r ^ {n - 1} - 1]},\\lambda T, T ^ {[r ^ n - r ^ {n - 1}]}}],\n\t\\end{equation*}\n\tand $\\lambda_{i(n) - r ^ {n - 1}}=\\lambda$.\n\tIn a similar way to the proof of Theorem \\ref{thm:main1}, we obtain\n\t\\begin{gather*}\n\t\\abs{\\Phi - \\beta_n} \\asymp q ^ {- 2 \\ell (r ^ n - 1) \/ (r - 1) - 2 m r ^ n},\\quad \n\t\\abs{\\beta_n - \\beta_n'} \\asymp q ^ {- 2 \\ell (r ^ n - 1) \/ (r - 1) + 2 r ^ {n - 1}},\\\\\n\tH(\\beta_n) \\asymp q ^{2 \\ell (r ^ n - 1) \/ (r - 1) + r ^ n - 2 r ^ {n - 1}}.\n\t\\end{gather*}\n\tTherefore, we have \\eqref{eq:mainequal6} and \\eqref{eq:mainequal66} by Lemma \\ref{lem:bestquad2}.\n\tIf \\eqref{eq:maincond7} holds, then we obtain \\eqref{eq:mainequal7} and \\eqref{eq:mainequal77} by Lemma \\ref{lem:bestquad}.\n\\end{proof}\n\n\n\n\\begin{proof}[Proof of Corollary \\ref{cor:main2}]\n\tIn the same way to the proof Corollary \\ref{cor:main1}, we prove Corollary \\ref{cor:main2}.\n\\end{proof}\n\n\n\n\n\n\\subsection*{Acknowledgements}\nThe author would like to thank the referee for helpful comments.\n\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\nIn the eighteenth century, the famous astronomer William Herschel showed us the\npowerful method of star count to understand our own Milky Way \\citep{Hers1785}.\nThe technique has been used since then by generations of astronomers. With\ngreat improvement on data collection over the years, more and more details of\nour Milky Way were unfolded. However, the characteristic scales of smooth\nGalactic structures (i.e., disks and halo) obtained by previous studies does\nnot converge to a common value as an outcome of improving data collection\n(see Table~\\ref{table_ref}). The spread of values is attributed to the\ndegeneracy of Galactic model parameters (i.e., same star count data could be\nfitted equally well by different Galactic models) \\citep{Chen2001,\nSiegel2002, Juric2008, Bilir2008}.\nThis is due to the different sky regions\nand limiting magnitudes (i.e., limiting volumes) used in these studies\n\\citep{Siegel2002, Karaali2004, Bilir2006a, Bilir2006b, Juric2008}. On the\ncontrary, \\citet{Bilir2008, Yaz2010} did not show such degeneracy in determining\nthe Galactic model parameters. This controversy is still\nactively debated.\nTherefore, systematic all sky surveys with deeper limiting magnitude and wider sky\nregion, such as the Two Micron All Sky Survey \\citep[2MASS;][]{Skrutskie2006},\nthe Sloan Digital Sky Survey \\citep[SDSS;][]{York2000}, the Panoramic Survey\nTelescope \\& Rapid Response System \\citep[Pan-Starrs;][]{Kaiser2002} and the\nGAIA mission \\citep{Perryman2001} could provide a good opportunity for us to\nstudy our Galaxy from a global\nperspective. Free from limited sky\nfields, astronomers can acquire many more information from the stellar distribution\nof these surveys.\n\nOn Galactic structure study, besides the simple and smooth\ntwo-component model \\citep{Bahcall1980} or the three-component model\n\\citep{Gilmore1985}, many more structures have been discovered, such as\ninner bars in the Galactic center\n\\citep{Alves2000, Hammersley2000, vanLoon2003, Nishiyama2005, Cabrera2008},\nand flares and warps \\citep{Lopez2002, Robin2003, Momany2006, Reyle2009},\nwhich has been contributed the variation of disk model parameters with\nGalactic longitude \\citep{Cabrera2007, Bilir2008, Yaz2010}.\nMoreover, the overdensities in the halo, such as Sagittarius\n\\citep{Majewski2003}, Triangulum-Andromeda \\citep{Rocha2004, Majewski2004},\nVirgo \\citep{Juric2008}, and in the outer disk, such as Canis Major\n\\citep{Martin2004}, Monoceros \\citep{Newberg2002,Rocha2003}, show the\ncomplexity of the Milky Way. The formation history of our Galaxy is more\ncomplicated than what we thought previously.\n\nOn stellar luminosity function\nstudy, a recent study by \\citet{Chang2010} using the $K_s$ band star count of\n2MASS point source catalog \\citep[2MASS PSC;][]{Cutri2003} verified for\nthe first time the universality hypothesis of the luminosity function (i.e.,\nthe common practice of assuming one luminosity function for the entire Milky\nWay).\n\nWe are interested in the global smooth structure of our Galaxy. In view of the\nexisting fine structures (e.g., flares, warps, overdensities\\dots etc.), the\nstructure of the smooth components can be determined either by including these\nfine structures in a grand Galaxy model or by avoiding the sky area\n``contaminated'' with these features. The first method demands a complex model,\nwhich involves many more structure parameters, and needs high computing power\nto accomplish the fitting task. \\citet{Lopez2002} is a good example of this\nmethod using 2MASS data. The second method is clearly simpler but needs\njustifications. We observe that (1) the fine structures (e.g., inner bars,\nflares and warps), which have observable contribution on 2MASS star count data,\nare all confined in the Galactic plane region.\n(2) The overdensities or substructures in the outer disk region and the halo are\ndifficult to identify in general.\nTheir contribution is negligible in the 2MASS star count data.\nHere we quote \\citet{Majewski2004} on halo substructure:\n``This substructure is typically subtle and obscured by a substantial\nforeground veil of disk stars, eliciting its presence requires strategies that\noptimize the substructure signal compared to the foreground noise.''\n\nWe prefer the second method in this paper and only use Galactic latitude\n$|b|>30^{\\circ}$ 2MASS $K_s$ band star\ncount data to obtain the structure parameters of a three-component model.\nIn addition, the influence of the near infrared extinction in these\nregions is small. This also allows us to use a simpler extinction model for\ncorrection. We describe our model in section 2 and the analysis method in\nsection 3. Section 4 provides the results and a discussion.\n\n\n\n\\section{The Milky Way Model}\nWe adopt a three-component model for the smooth stellar distribution of the\nMilky Way. It comprises a thin disk, a thick disk and an oblate halo\n\\citep{Bahcall1980, Gilmore1983}. The total stellar density $n(R,Z)$\nat a location $(R,Z)$\nis the sum of the thin disk $D_1$, the think disk $D_2$ and the\nhalo $H$,\n\\begin{equation}\\label{density}\n n(R,Z)=n_0\\left[D_1(R,Z)+D_2(R,Z)+H(R,Z)\\right]\\,,\n\\end{equation}\nwhere $R$ is the galactocentric distance on the Galactic plane, $Z$ is the\ndistance from the the Galactic mid-plane and $n_0$ is the local stellar\ndensity of the thin disk at the solar neighborhood.\n\nThe stellar distribution of the thin disk $D_1$ and the thick disk $D_2$\ndecreases exponentially along $R$ and $Z$ (the so called double exponential\ndisk),\n\\begin{equation}\\label{Di}\n D_i(R,Z)=f_i\\exp\\left[-\\,{(R-R_\\odot)\\over H_{ri}}-\\,{(|Z|-|Z_\\odot|)\\over H_{zi}}\\right]\\,,\n\\end{equation}\nwhere $(R_\\odot,Z_\\odot)$ is the location of the Sun, $H_{ri}$ is the\nscale-length, $H_{zi}$ is the scale-height, and $f_i$ is the density ratio to\nthe thin disk at the solar neighborhood. The subscript $i=1$ stands for the\nthin disk (thus $f_1=1$) and $i=2$ for thick disk. We adopted $R_\\odot=8$ kpc\nin our model \\citep{Reid1993}.\n\nThe halo is a power law decay oblate spheroid flattening in the $Z$ direction,\n\\begin{equation}\\label{SH}\n H(R,Z)=f_h\\left[R^2+(Z\/\\kappa)^2\\over R_\\odot^2+(Z_\\odot\/\\kappa)^2\\right]^{-p\/2}\\,,\n\\end{equation}\nwhere $\\kappa$ is the axis ratio, $p$ is the power index and $f_h$ is the local\nhalo-to-thin disk density ratio.\n\n\\citet{Chang2010} showed that the whole sky $K_s$ band luminosity function can\nbe well approximated by a single power law with a power law index\n$\\gamma=1.85\\pm0.035$, a bright cutoff at $M_b=-7.86\\pm0.60$ and a faint cutoff\nat $M_f=6.88\\pm0.66$. We adopt this luminosity function in the following analysis.\nIn normalized form, it is\n\\begin{equation}\\label{LF}\n \\psi(M_{K_s})={2\\log_{\\rm e}10\\ (\\gamma-1)\\over\n 5\\left[10^{2(\\gamma-1)M_{f}\/5}-10^{2(\\gamma-1)M_{b}\/5}\\right]}\\,10^{2(\\gamma-1)M\/5}\\,.\n\\end{equation}\nNote that $\\psi(M)$ includes all luminosity classes.\nIn our analysis the observing magnitude range is $5 \\le K_s \\le 14$ mag.\nThe corresponding distances\nof bright cutoff, faint cutoff and $M_{K_s}=0$ are 3.4 kpc to 216 kpc, 4 pc to\n265 pc and 0.1 kpc to 6.3 kpc, respectively.\n\nAlthough we only use NIR data in the medium and high Galactic latitude regions,\nwe still need to correct possible interstellar extinction. We adopt the new\nCOBE\/IRAS result \\citep{Chen1999} and convert it to the $K_s$ band extinction\nby $A_{K_s}\/E(B-V)=0.367$ \\citep{Schlegel1998}. This extinction model is then\napplied to our simulation data. The extinction values of most of our analyzed\nregions are $A_{K_s}< 0.03$.\n\n\n\n\\section{The Data and Analysis Method}\n\\subsection{The 2MASS Data}\nThe 2MASS Point Source Catalog \\citep[2MASS PSC,][]{Cutri2003} is employed to\ncarry out the $K_s$ band differential star count for the entire Milky Way. We\ndivide the whole sky into 8192 nodes according to level 5 Hierarchical\nTriangular Mesh \\citep[HTM,][]{Kunszt2001}. The level 5 HTM samples the whole\nsky in roughly equal area with an average angular distance about 2 degrees\nbetween any two neighboring nodes. The amount of stars within 1 degree radius\nof each node (i.e., each node covers $\\pi$ square degree) is then retrieved with\na bin size $K_s$=0.5 mag via 2MASS online data service\n\\citep[][catalog]{Cutri2003}. Our selection criterion is: the object must be\ndetected in all $J, H, K_s$ bands and has signal-to-noise ratio $\\ge 5$. Since\nthe limiting magnitude of 2MASS $K_s$ band is 14.3 mag, which has 10\nsignal-to-noise ratio and 99\\% completeness \\citep[see Table~1\nin][]{Skrutskie2006}, and $K_s \\le 5$ mag objects have relatively large\nphotometric error, we only compare 2MASS data with our simulation data from\n$K_s=5$ to $14$ mag.\n\nIn order to minimize the effects coming from the close-to-Galactic-plane fine\nstructures (e.g., flares, warps, arms, budge and bars\\dots etc., which have\nconsiderable contribution to the star count data), and the relatively complex\nextinction correction at the low Galactic latitude region,\nwe avoid low galactic latitudes, and only consider data in Galactic\nlatitude $|b|>30^{\\circ}$. Although several\noverdensities in the halo (such as Sagittarius, Triangulum-Andromeda,\nVirgo\\dots etc.) were identified, they cannot be\npicked up\nfrom the\noverwhelming foreground field stars on star count data without additional\ninformation \\citep[e.g., color, distance, metallicity\\dots\netc.,][]{Majewski2004}. Therefore, their contribution to the 2MASS $K_s$ band\ndifferential star count is negligible and will not affect our result. We also\nexclude the areas around Large and Small Magellanic Clouds for their\nsignificant stellar population.\n\n\n\\subsection{The Analysis Method}\\label{method}\nThe Maximum Likelihood Method (Bienayme et al. 1987) is applied to compare the\n$K_s$ band 2MASS differential star counts and the simulation data to search for\nthe best-fit structure parameters of the three-component Milky Way model. Our\nfitting strategy is as follows:\n\\begin{enumerate}\\label{fitting_steps}\n\\item Take $R_\\odot=8$ kpc;\n\\item choose one $Z_\\odot$ and work out the maximum likelihood value by fitting the 9 parameters\n$(n_0,H_{z1},H_{r1},f_2,H_{z2},H_{r2},f_h,\\kappa,p)$;\n\\item repeat step 2 for other $Z_\\odot$;\n\\item pick the $Z_\\odot$ corresponds to the maximum of the maximum likelihood values in step 3;\n\\item repeat steps 2 to 4 for finer grid size of the 9-parameter fit and a narrower range of\n$Z_\\odot$ around the one found in step 4;\n\\item the uncertainty is estimated by adding Poisson noise on the simulation data to see\nhow the likelihood varies. The difference of the likelihoods of 500\nrealizations of the same model, differed by the Poisson statistics only,\ngives a range of likelihood around the maximum likelihood that defines the\nconfidence level.\n\\end{enumerate}\nTable~\\ref{para_space} lists our searching parameter space and the finest grid\nsize we used. The key parameters in our study are $n_0$, $H_{z1}$, $f_2$ and\n$H_{z2}$ (see Eqs. (1)-(2)). The first two play a primary role on the variation\nof the likelihood value and the latter two play a secondary role. The other\nfive parameters $H_{r1}$, $H_{r2}$, $f_h$, $\\kappa$ and $p$ (see Eqs. (1)-(3))\nare non-key parameters, which play a minor role and do not affect the\nlikelihood value as much as the key parameters.\n\n\\section{The Results and Discussion}\nTable~\\ref{result} lists our best-fit results with the corresponding\nuncertainties. Fig.~\\ref{2Dcon} shows contour plots of likelihood against\ndifferent pairs of parameters. The contour changes dramatically along the key\nparameters $n_0$, $H_{z1}$, $f_2$ and $H_{z2}$, but relatively mild along other\nparameters $H_{r1}$, $H_{r2}$, $f_h$, $p$ and $\\kappa$.\nThis indicates the importance of the key parameters in determining the best-fit result.\n\nDegeneracies exist between some pairs of key parameters, such as $(n_0,\nH_{z1})$, $(n_0, f_2)$, $(H_{z1}, f_2)$\\dots etc. Here degeneracy means that\nthe likelihood value stays almost the same when the pairs of parameters change\ntogether in a particular way. Similar degeneracy between the local\nthick-to-thin disk ratio $f_2$ and the scale-height of the thick disk $H_{z2}$\nhas been reported in \\citet{Chen2001,Siegel2002,Juric2008,Bilir2008}.\nConsequently, it is possible that different combinations of parameters can be\nregarded as `acceptable' fitting. For example, if we choose a higher local\nstellar density $n_0$, then we can pick a smaller $H_{z1}$ such that the\nlikelihood value is very close to the maximum likelihood value and assign it as\nthe `best-match' scale-height of the thin disk. Therefore, a thin light color\ndiagonal strip shows on the $n_0$ against $H_{z1}$ contour plot in\nFig.~\\ref{2Dcon}. Besides, similar trends happen in other pairs of key\nparameters. When one parameter of the pair is higher, we can get a similar\nlikelihood value by lowering the other parameter (see the corresponding\ntwo-parameter contour plots of key parameters in Fig.~\\ref{2Dcon}). This\nanti-correlation is not unexpected. The number of stars along the line of sight\nin the model increases when any one of the key parameters increases. Thus for a\ngiven observed number of stars, an increase in one key parameter can be\ncompensated by a decrease in the other. Perhaps this is the reason that our\nbest-fit scale-height of the thin disk, $H_{z1}=360$ pc, is somewhat larger\nthan the reported values, $H_{z1}=285$ pc, in \\citet{Lopez2002} study (they\nalso use star count of 2MASS to obtain Galactic model parameters). Our best-fit\nlocal stellar density from $K_s=$-8 to 6.5 mag is $n_0\\sim 0.030$ star\/pc$^3$,\nwhich is about half of $\\sim 0.056$ star\/pc$^3$ of the corresponding value\ncited in \\citet{Eaton1984}, and a bit lower than $\\sim 0.032$ star\/pc$^3$ of\nthe corresponding value cited in \\citet{Lopez2002}. As a result, our fitting\ntends to choose a larger scale-height of the thin disk. If we force our local\nstellar density to be comparable with that of \\citet{Lopez2002}, then the\ncorresponding `best-fit' scale-height of the thin disk would be $\\sim$320 pc,\nwhich is closer to their result. If we choose even higher local stellar\ndensity, then the `best-fit' scale-height of the thin disk would be made\nbetween 200 to 300 pc, which is similar to the most of recent studies (see\nTable~\\ref{table_ref}). Moreover, we do not apply binarism correction in our\nanalysis, and it has been shown that scale-length and scale-height might be\nunderestimated without binarism correction \\citep{Siegel2002, Juric2008,\nIvezic2008, Yaz2010}.\n\n\nFor our purpose, we deem that in order to lift the degeneracy it is crucial to\nhave a reliable near infrared luminosity function by observation or a near\ninfrared local stellar density. Unfortunately, a systematic study in this\ndirection is yet to come. Some related studies, such as synthetic luminosity\nfunction \\citep[see e.g.,][]{Girardi2005} or luminosity function transformed\nfrom optical observation \\citep[see e.g.,][]{Wainscoat1992} do exist, but some\nuncertainties still need to be settled (e.g., the initial mass function,\nmass-luminosity relation for NIR and color transformation between different\nwavelengths\\dots etc.). Once the `true' local stellar density is known, the\n`true' structure of the Milky Way would be revealed.\n\n\n\nIn order to see how good the agreement between 2MASS data and our best-fit\nmodel, we show an all sky map of the ratios of observed to predicted integrated\nstar count from $K_s=$5 to 14 mag for each node as a function of position on\nthe sky in Fig.~\\ref{comparison}, and the color indicates the values of the\nratios. We do not see obvious deviation in the Galactic latitude\n$|b|>30^{\\circ}$ areas, but only in the Large Magellanic Cloud and the Small\nMagellanic Cloud areas. For comfirmation, we plot integrated star count from\n$K_s=$5 to 14 mag along Galactic longitude and latitude for $|b|>30^{\\circ}$\nareas in Figs.~\\ref{surf_b} \\& \\ref{surf_l}. We see 2MASS data and our best-fit\nmodel agree well and only some small deviations in the nodes at the\nanti-Galactic center $b\\sim30^{\\circ}$ areas. The significant spikes in\nFigs.~\\ref{surf_b} \\& \\ref{surf_l} are due to the populations of the Large\nMagellanic Cloud and the Small Magellanic Cloud. For testing how these small\ndeviations affect the key parameters selection, we exclude these small\ndeviations and re-analyze the data with fixed non-key parameters. It shows only\nthe scale-height of thick disk shifts slightly but still within our error\nestimation. Besides, we do not see any significant differences of reported\noverdensities in sky regions with $|b|>30^{\\circ}$. Thus, we conclude that the\nthree-component model can describe the Milky Way structure sufficiently well\nfor high Galactic latitude regions and the single power law luminosity function\nof \\citet{Chang2010} is a good approximation as well. Since our main purpose is\nto search for the global Galactic model parameters, we do not try to explore\nthe best-fit result for each node individually as what \\citet{Cabrera2007,\nBilir2008, Yaz2010} have done in their studies to seek the variations of disk\nmodel parameters with Galactic longitude. Instead, we treat the differences of\n2MASS data to our best-fit model as deviations from a global smooth\ndistribution. We believe that similar variations in disk model parameters will\nbe obtained if we take similar analysis procedure (i.e., searching best-fit\nparameters for each node), but this is beyond the scope of this work.\nBecause we do not consider flares, warps and other overdensities in our\nmodel, there are some discrepancies between 2MASS data and the model in the low\nGalactic latitude regions (see Fig.~\\ref{comparison}). These fine structures\nmake the star distribution more fluffy in the vertical direction toward the\nedge of the Milky Way. Hence the discrepancy between 2MASS data and the model\nincreases vertically towards the anti-Galactic center region. In addition, the\nabsence of Galactic bulge in our model contributes to the large discrepancy in\nGalactic center areas. The difference in the low Galactic latitude region needs\nmore delicate analysis to rectify \\citep[see, e.g.,][]{Lopez2002, Momany2006}.\n\n\\subsection{Summary}\nIn summary, we set forth to study the global smooth structure of the Milky Way\nby a three-component stellar distribution model which comprises two double\nexponential disks (one thin and one thick) and an oblate halo. The $K_s$ band\n2MASS star count is used to determine the structure parameters. To avoid the\ncomplication introduced by the fine structures and complex extinction\ncorrection close to Galactic plane, we use only Galactic latitude\n$|b|>30^{\\circ}$ data. There are 10 parameters in the model, but only four of\nthem play the dominant role in the fitting process. They are the local stellar\ndensity of the thin disk $n_0$, the local density ratio of thick-to-thin disk\n$f_2$, and the scale-height of the thin and thick disks $H_{z1}$ and $H_{z2}$.\nThe best-fit result is listed in Table~\\ref{result}. In short the scale-height\nof the thin and thick disks are $360\\pm 10$ pc and $1020\\pm 30$ pc,\nrespectively; the scale-length of the thin disk is $3.7\\pm 1.0$ kpc and that of\nthick disk is $5.0\\pm 1.0$ kpc (the uncertainty in scale-length is large\nbecause it is not very sensitive to high latitude data.) The local stellar\ndensity ratio of thick-to-thin disk and halo-to-thin disk are $7\\pm 1$\\% and\n$0.20\\pm 0.10$\\%, respectively. The local stellar density of the thin disk is\n$0.030\\pm 0.002$ stars\/pc$^3$.\n\nAn all sky comparison of the 2MASS data to our best-fit model is shown in\nFig.~\\ref{comparison}. A good agreement in the Galactic latitude\n$|b|>30^{\\circ}$ areas is expected from our fitting procedure. In low Galactic\nlatitude regions, fine structures (such as flares, warps\\dots etc.) increase\nthe effective scale-height towards the edge of the Milky way. This is reflected\nin the fan-like increase in discrepancy towards the anti-Galactic center\nregions.\n\nDegeneracy (i.e., different combinations of parameters give similar likelihood\nvalues) is found in pairs of key parameters\n(see Fig.~\\ref{2Dcon}). Thus different combinations of parameters may fit the\ndata almost as good as the best-fit one, and these are all legitimate\n`acceptable' fitting in view of the uncertainty. Therefore, accompanying our\nlower local stellar density $0.030$ stars\/pc$^3$ (from $K_s=-$8 to 6.5 mag) is\na higher thin disk scale-height $360$ pc. In the context of NIR star count, the\nNIR luminosity function or the NIR local stellar density is imperative to\ndetermine the scale-height and other Milky Way structure parameters. We hope\nthat systematic study on the luminosity function and the local stellar density\nin near infrared will be available in the near future.\n\n\n\\acknowledgements\nWe acknowledge the use of the Two Micron All Sky Survey Point\nSource Catalog (2MASS PSC).\nWe would like to thank the anonymous referee, whose advice greatly improves the\npaper.\nCMK is grateful to S. Kwok and K.S. Cheng for their\nhospitality during his stay at the Department of Physics, Faculty of Science,\nUniversity of Hong Kong. This work is supported in part by the National Science\nCouncil of Taiwan under the grants NSC-98-2923-M-008-001-MY3 and\nNSC-99-2112-M-008-015-MY3.\n\n\n\\begin{deluxetable}{lllllllll}\n\\tabletypesize{\\tiny} \\setlength{\\tabcolsep}{0.03in} \\tablecaption{Previous\nGalactic models. The parentheses are the corrected values for binarism. The\nasterisk denotes the power-law index replacing Re. References indicating with\nthe original result table mean Galactic longitude or limiting magnitude\ndependent Galactic model parameters. \\label{table_ref}} \\tablewidth{0pt}\n\\tablehead{ \\colhead{$H_{z1}$ (pc)} & \\colhead{$H_{r1}$ (kpc)} & \\colhead{$f_2$\n(\\%)} & \\colhead{$H_{z2}$ (kpc)} & \\colhead{$H_{r2}$ (kpc)} &\n\\colhead{$f_h$(\\%)} & \\colhead{Re(S) (kpc)} & \\colhead{$\\kappa$} &\n\\colhead{Reference}} \\startdata\n310 - 325 & - & 0.0125-0.025 & 1.92 - 2.39 & - & - & - & - & \\citet{Yoshii1982} \\\\\n300 & - & 0.02 & 1.45 & - & - & - & - & \\citet{Gilmore1983} \\\\\n325 & - & 0.02 & 1.3 & - & 0.002 & 3 & 0.85 & \\citet{Gilmore1984} \\\\\n280 & - & 0.0028 & 1.9 & - & 0.0012 & - & - & \\citet{Tritton1984} \\\\\n125-475 & - & 0.016 & 1.18 - 2.21 & - & 0.0013 & 3.1* & 0.8 & \\citet{Robin1986} \\\\\n300 & - & 0.02 & 1 & - & 0.001 & - & 0.85 & \\citet{delRio1987} \\\\\n285 & - & 0.015 & 1.3 - 1.5 & - & 0.002 & 2.36 & Flat & \\citet{Fenkart1987} \\\\\n325 & - & 0.0224 & 0.95 & - & 0.001 & 2.9 & 0.9 & \\citet{Yoshii1987} \\\\\n249 & - & 0.041 & 1 & - & 0.002 & 3 & 0.85 & \\citet{Kuijken1989} \\\\\n350 & 3.8 & 0.019 & 0.9 & 3.8 & 0.0011 & 2.7 & 0.84 & \\citet{Yamagata1992} \\\\\n290 & - & - & 0.86 & - & - & 4 & - & \\citet{vonHippel1993} \\\\\n325 & - & 0.020-0.025 & 1.6-1.4 & - & 0.0015 & 2.67 & 0.8 & \\citet{ReidMajewski1993} \\\\\n325 & 3.2 & 0.019 & 0.98 & 4.3 & 0.0024 & 3.3 & 0.48 & \\citet{Larsen1996} \\\\\n250-270 & 2.5 & 0.056 & 0.76 & 2.8 & 0.0015 & 2.44 - 2.75* & 0.60 - 0.85 & \\citet{Robin1996, Robin2000} \\\\\n260 & 2.3 & 0.074 & 0.76 & 3 & - & - & - & \\citet{Ojha1996} \\\\\n290 & 4 & 0.059 & 0.91 & 3 & 0.0005 & 2.69 & 0.84 & \\citet{Buser1998, Buser1999} \\\\\n240 & - & 0.061 & 0.79 & - & - & - & - & \\citet{Ojha1999} \\\\\n280\/267 & - & 0.02 & 1.26\/1.29 & - & - & 2.99* & 0.63 & \\citet{Phleps2000} \\\\\n330 & 2.25 & 0.065 - 0.13 & 0.58 - 0.75 & 3.5 & 0.0013 & - & 0.55 & \\citet{Chen2001} \\\\\n- & 2.8 & 3.5 & 0.86 & 3.7 & - & - & - & \\citet{Ojha2001} \\\\\n280(350) & 2 - 2.5 & 0.06 - 0.10 & 0.7 - 1.0 (0.9 - 1.2) & 3 - 4 & 0.0015 & - & 0.50 - 0.70 & \\citet{Siegel2002} \\\\\n285 & 1.97 & - & - & - & - & - & - & \\citet[][]{Lopez2002} \\\\\n - & 3.5 & 0.02-0.03 & 0.9 & 4.7 & 0.002-0.003 & 4.3 & 0.5-0.6 & \\citet{Larsen2003} \\\\\n320 & - & 0.07 & 0.64 & - & 0.00125 & - & 0.6 & \\citet{Du2003} \\\\\n265-495 & - & 0.052-0.098 & 0.805-0.970 & - & 0.0002-0.0015 & - & 0.6-0.8 & \\citet[][Table 16]{Karaali2004} \\\\\n268 & 2.1 & 0.11 & 1.06 & 3.04 & - & - & - & \\citet{Cabrera2005} \\\\\n300 & - & 0.04-0.10 & 0.9 & - & - & 3\/2.5* & 1\/0.6 & \\citet{Phleps2005} \\\\\n220 & 1.9 & - & - & - & - & - & - & \\citet{Bilir2006a} \\\\\n160-360 & - & 0.033-0.076 & 0.84-0.87 & - & 0.0004-0.0006 & - & 0.06-0.08 & \\citet[][Table 15]{Bilir2006b} \\\\\n301\/259 & - & 0.087\/0.055 & 0.58\/0.93 & - & 0.001 & - & 0.74 & \\citet[][Table 5]{Bilir2006c} \\\\\n220-320 & - & 0.01-0.07 & 0.6-1.1 & - & 0.00125 & - & $>$0.4 & \\citet{Du2006} \\\\\n206\/198 & - & 0.16\/0.10 & 0.49\/0.58 & - & - & - & 0.45 & \\citet{Ak2007} \\\\\n140-269 & - & 0.062-0.145 & 0.80-1.16 & - & - & - & - & \\citet[][Table 1]{Cabrera2007} \\\\\n220-360 & 1.65-2.52 & 0.027-0.099 & 0.62-1.03 & 2.3-4.0 & 0.0001-0.0022 & - & 0.25-0.85 & \\citet[][Table 3]{Karaali2007} \\\\\n167-200 & - & 0.055-0.151 & 0.55-0.72 & - & 0.0007-0.0019 & - & 0.53-0.76 & \\citet[][Table 1]{Bilir2008} \\\\\n245(300) & 2.15(2.6) & 0.13(0.12) & 0.743(0.900) & 3.261(3.600) & 0.0051 & 2.77* & 0.64 & \\citet{Juric2008} \\\\\n325-369 & 1.00-1.68 & 0.0640-0.0659 & 0.860-0.952 & 2.65-5.49 & 0.0033-0.0039 & - & 0.0489-0.0654 & \\citet[][Table 1]{Yaz2010} \\\\\n360 & 3.7 & 0.07 & 1.02 & 5 & 0.002 & 2.6* & 0.55 & This Work \\\\\n\\enddata\n\\end{deluxetable}\n\n\n\\begin{deluxetable}{lcc}\n\\tabletypesize{\\scriptsize} \\tablecaption{The searching parameter space.\n\\label{para_space}} \\tablewidth{0pt} \\tablehead{ \\colhead{} & \\colhead{Range} &\n\\colhead{Grid size}} \\startdata\nThin Disk & \\\\\n\\quad $H_{r1}$ & 1.0-6.0 kpc & 100 pc \\\\\n\\quad $H_{z1}$ & 200-450 pc & 10 pc \\\\\n\\quad $n_0$ & 0.02-0.04 stars\/pc$^3$ & 0.002 stars\/pc$^3$ \\\\\n\\quad $Z_\\odot$ & 11-35 pc & 2.5 pc \\\\\nThick Disk & \\\\\n\\quad $H_{r2}$ & 3.0-7.0 kpc & 200 pc \\\\\n\\quad $H_{z2}$ & 400-1200 pc & 20 pc \\\\\n\\quad $f_2$ & 0-20\\% & 1\\% \\\\\nSpheroid & \\\\\n\\quad $\\kappa$ & 0.1-1.0 & 0.05 \\\\\n\\quad $p$ & 2.3-3.3 & 0.1 \\\\\n\\quad $f_h$ & 0-0.45\\% & 0.05\\% \\\\\n\\enddata\n\\end{deluxetable}\n\n\\begin{deluxetable}{lcc}\n\\tabletypesize{\\scriptsize} \\tablecaption{The best-fit Milky Way model.\n\\label{result}} \\tablewidth{0pt} \\tablehead{ \\colhead{} & \\colhead{Value} &\n\\colhead{Uncertainty}} \\startdata\nThin Disk & \\\\\n\\quad $H_{r1}$ & 3.7 kpc & 1.0 kpc \\\\\n\\quad $H_{z1}$ & 360 pc & 10 pc \\\\\n\\quad $n_0$ & 0.030 stars\/pc$^3$ & 0.002 stars\/pc$^3$ \\\\\n\\quad $Z_\\odot$ & 25 pc & 5 pc \\\\\nThick Disk & \\\\\n\\quad $H_{r2}$ & 5.0 kpc & 1.0 kpc \\\\\n\\quad $H_{z2}$ & 1020 pc & 30 pc \\\\\n\\quad $f_2$ & 7\\% & 1 \\% \\\\\nSpheroid & \\\\\n\\quad $\\kappa$ & 0.55 & 0.15 \\\\\n\\quad $p$ & 2.6 & 0.6 \\\\\n\\quad $f_h$ & 0.20\\% & 0.10\\% \\\\\n\\enddata\n\\end{deluxetable}\n\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzziujm b/data_all_eng_slimpj/shuffled/split2/finalzziujm new file mode 100644 index 0000000000000000000000000000000000000000..5f782669f783310d52f9a2e4fe683fede0d806a5 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzziujm @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\n\\label{section-introduction}\nGiven positive integers $m$ and $n$ with $m - n\\equiv 0 \\bmod 8$, let\n$\\mathrm{II}_{m,n}$ denote the unique even self-dual lattice of signature\n$(m,n)$; see \\cite{CS:SPLAG} chapter 26, 27.\nFor basic definitions about lattices, see\n\\cite{REB:L} or \\cite{CS:SPLAG}. All the lattices considered here are\nnonsingular unless otherwise stated. A bilinear form is usually denoted by $\\ip{\\;}{\\;}$.\nThe lattice $\\mathrm{II}_{1,1}$ is sometimes called a {\\it hyperbolic cell}.\nThe even self-dual Lorentzian lattices are $\\mathrm{II}_{8n + 1, 1}$. Among these,\nthe lattice $L = \\mathrm{II}_{25,1}$ stands out\nas exceptional (see \\cite{CS:SPLAG} or \\cite{REB:Thesis}, \\cite{REB:L}, \nfor a wealth of information on $L$).\nOf course this is related to the many exceptional properties of the\nLeech lattice $\\Lambda$. \nIndeed $ L \\simeq \\Lambda \\oplus \\mathrm{II}_{1,1}$.\nLet $L(2)$ denote the set of norm $2$ vectors (also called {\\it roots}) of $L$.\nThe roots form a single orbit under the automorphism group $\\op{Aut}(L)$.\nIn this article, we study the Poincare-Weierstrass series\n\\begin{equation*}\nE_L(z,s) = E(z,s) = \\sum_{r \\in L(2) } \\ip{r}{z}^{-s}.\n\\end{equation*}\nHere $z$ is a negative norm vector in $L \\otimes \\mathbb{R}$ and\n $s$ is a complex\nnumber.\nThe function $E(z,s)$ is real analytic in $z$. \nComplex analytic versions of this function\nhave appeared in Looijenga's study of compactifications of\ncomplex ball quotients (see the functions $F^{(l)}_{\\mathcal{O}}$ in lemma 5.4 of \\cite{L:CDBA1}). \nThe functions $F^{(l)}_{\\mathcal{O}}$ are analogous to the \nEisenstein series $\\sum_{(m,n) \\neq (0,0)} (m \\tau + n)^{-2l}$ whereas\n $E(z,s)$ is analogous to real analytic Eisenstein series.\n\\par\nThe infinite series for $E(z,s)$ converges for $\\op{Re}(s) > 25$\n(For a conceptual proof, see lemma 5.4 of \\cite{L:CDBA1}.\nFor the sake of completeness, we have included an elementary\nproof in appendix \\ref{l-convergence}).\nSo, for fixed $s$ with $\\op{Re}(s) > 25$, we obtain a real analytic function\n$z \\mapsto E(z,s)$ on hyperbolic space\n$\\mathcal{H}^{25}$.\nBy definition, $E(z,s)$ is an automorphic function invariant under $\\op{Aut}(L)$ with singularities\nalong the mirrors of the reflection group of $L$. \nOur objective is to compute the Fourier series for $E(z,s)$ and prove that\n$E(z,s)$ can be analytically continued to a meromorphic function on $\\op{Re}(s) > 25\/2$.\n\\par\nMuch of what we say about $E_{\\mathrm{II}_{25,1}}(z,s)$\ncan probably be generalized for all the even self-dual \nLorentzian lattices $\\mathrm{II}_{ 8 n + 1, 1}$\n(at least if we assume $n \\geq 3$, which would guarantee the existence of\na positive definite even self-dual lattice of rank $8n$ with no roots). \nHowever, in section \\ref{section-II251} and \\ref{section-Fourier-expansion}, while\nstudying $E_L(z,s)$, \nwe have decided to restrict to the example $L = \\mathrm{II}_{25,1}$ for several reasons.\nFirst, it allows us to keep the exposition relatively simple.\nSecond, in view of moonshine, this is the most important example. \nThird, because of special properties of the Leech lattice, in particular, \nbecause of Conway's beautiful description of the automorphism group of $L$, \nthis is also the nicest example. \n\\par\nTo describe the form of the Fourier expansion of $E(z,s)$, we need to first\nbriefly recall Conway's description of the {\\it reflection group} $R(L)$. By definition,\n$R(L)$ is the subgroup of $\\op{Aut}(L)$ generated by reflections in the roots of $L$.\nBoth $\\op{Aut}(L)$ and $R(L)$ act on\nthe hyperbolic space\n $ B(L) = \\lbrace x \\in L \\otimes \\mathbb{R} \\colon x^2 < 0 \\rbrace\/ \\mathbb{R}^*\n \\simeq \\mathcal{H}^{25}$.\nTo describe a fundamental domain of $R(L)$ acting on $\\mathcal{H}^{25}$,\nchoose a {\\it Leech cusp} $\\rho$: this means that $\\rho$ is a primitive \nnorm zero vector of $L$ and $\\rho^{\\bot}\/\\rho \\simeq \\Lambda$.\nThis lets us split a hyperbolic cell from $L$ and \nidentify $L = \\Lambda \\oplus \\mathrm{II}_{1,1}$. So we write $v \\in L$ in the form\n$v = (\\lambda; m,n)$ with $\\lambda \\in \\Lambda$ and $m,n \\in \\mathbb{Z}$\nwith $v^2 = \\lambda^2 - 2 m n$. In this co-ordinate system $\\rho = (0; 0, 1)$.\nThe Leech cusp $\\rho$ determines a point in $\\partial \\mathcal{H}^{25}$,\nalso denoted by $\\rho$. The action of $R(L)$ on $\\mathcal{H}^{25}$\nhas a unique fundamental domain $C$ whose closure contains $\\rho$.\nWe say that $C$ is the {\\it Weyl chamber ``around\" $\\rho$}.\nThe walls of $C$ are in bijection with the vectors of Leech lattice. The angles\nbetween these walls are determined by the inner products of the corresponding Leech lattice\nvectors. This lets one describe $R(L)$ explicitly as a Coxeter group so\nthat\n{\\it ``the Leech lattice is the Dynkin diagram for $R(L)$\"}.\n(see \\ref{t-C} or \\cite{C:Aut26}, \\cite{REB:Thesis} for details).\n\\par\nBecause the Leech lattice has no roots, there are no mirrors of $R(L)$\nthat pass through $\\rho$. The subgroup $\\mathbb{T}$\nof $\\op{Aut}(L)$ that fix $\\rho$ and act trivially on\n$\\rho^{\\bot}\/\\rho$ is a free abelian group isomorphic to $\\Lambda$.\nWe call $\\mathbb{T}$ the group of {\\it translations}. \nWe compute the Fourier series of $E(z,s)$ at $\\rho$\nwith respect to the group of translations; see \\ref{topic-continuous-group-of-translations},\n\\ref{topic-discrete-group-of-translations}, \\ref{topic-setup-for-Fourier-series-computation}\n for details.\nFix positive real numbers $h,k$ such that $2h^2\/k < 1$. Let $z \\in L \\otimes \\mathbb{R}$ \nbe a vector of norm $-k$ having the form $z = (v h; h, *)$, where\n$v \\in \\Lambda \\otimes \\mathbb{R}$ and the last coordinate of $z$ is determined by the condition\n$z^2 = -k$. Then $z$ (or rather its image in $\\mathcal{H}^{25}$)\nis contained in the Weyl chamber $C$. So $z$ does not lie on any mirror.\nThe Fourier series of $E(z,s)$ has the form\n\\begin{equation*}\nE((vh; h, *),s) = \\sum_{\\lambda \\in \\Lambda} a_{\\lambda}(k,h,s)\\exp( 2 \\pi i \\ip{\\lambda}{v}).\n\\end{equation*}\nEach coefficient in turn is an infinite series of the form\n\\begin{equation*}\na_{\\lambda}(k,h,s) = \\sum_{n = 1}^{\\infty} j_{\\lambda, n} n^{-s} g_{\\lambda}(s,1\/n)\n\\end{equation*}\nwhere \n\\begin{equation*}\nj_{\\lambda, n} = \\sum_{l \\in \\Lambda\/n \\Lambda : l^2\/2 \\equiv 1 \\bmod n} \n\\exp( 2 \\pi i \\ip{l}{\\lambda}\/n) \n\\end{equation*}\nand $g_{\\lambda}(s,y)$ is some function that is \nentire in $s$ and analytic for $y$ in a disc of radius bigger than $1$\nand exponentially decaying as $\\abs{\\lambda} \\to \\infty$; see theorem\n\\ref{th-fourier-coeffients}.\n\\par\nSection \\ref{section-j} is devoted to studying exponential sums of the form $j_{\\lambda, n}$\nfor general even self-dual lattices of arbitrary signature. \nThe results of section \\ref{section-j} and related results in appendix \\ref{appendix-gauss-sums}\nare probably known to experts. But since we could not find a reference,\nand we think these results may have applications in other contexts,\nwe have decided to include the details. In theorem \\ref{th-j} we compute\nthe sums $j_{\\lambda,n}$ in terms of Kloosterman sums.\nWeil's bound on Kloosterman sums implies that the infinite series for $a_{\\lambda}(k,h,s)$\nconverges for $\\op{Re}(s) > 25\/2$\nand it follows that $E(z,s)$ can be meromorphically continued in this half plane; \nsee theorem \\ref{th-analytic-continuation-of-E}.\nFurther, using analytic continuation of Kloosterman sum zeta functions,\nwe prove that the Dirichlet series $\\sum_{n} j_{\\lambda,n} n^{-s}$ and\nthe individual Fourier coefficients $a_{\\lambda}(k,h,s)$ have\nmeromorphic continuation to the whole $s$-plane \n(see \\ref{l-analytic-continuation-of-Dirichlet-series-of-j} and\n \\ref{remark-analytic-continuation-of-Fourier-coefficients}).\nHowever, since we are unable to find bounds on these functions obtained after\nanalytic continuation, we cannot prove continuation of $E(z,s)$ beyond\nthe line $\\op{Re}(s) = 25\/2$.\n\\par\nWe have two motivations for studying $E(z,s)$. \nThe first motivation is to find out if the residue of $E(z,s)$ at the ``first pole\"\nis related to some of the automorphic forms of type $O(25,1)$ constructed by taking Borcherds lift\n(see \\cite{REB:Aut}, section 10, in particular, the set up of example 10.7).\nThe second motivation is that the ``complex analytic versions of $E(z,s)$\" are\nrelevant in projective uniformization of a $13$ dimensional complex ball quotient\nwhose fundamental group is conjectured to be related to the monster simple group \\cite{Allcock-monstrous}.\nThese complex analytic versions of $E(z,s)$ are instances of the functions considered in \\cite{L:CDBA1}\nand they are defined as\n\\begin{equation*}\nE^{(l)}(z) = \\sum \\nolimits_{r \\in L(3) } \\ip{z}{r}^{-l}\n\\end{equation*}\nwhere $L$ is the unique\nhermitian $\\mathbb{Z}[e^{2 \\pi i\/3}]$-lattice of signature $(13,1)$ satisfying $\\sqrt{-3} L^{\\vee} = L$,\nthe sum is over $L(3)$ which is the set of roots of $L$, the exponent\n$l \\in 6 \\mathbb{N}$ and $z$ is a negative norm\nvector of $L \\otimes_{\\mathbb{Z}[e^{2 \\pi i\/3}]} \\mathbb{C}$. \nSo $E^{(l)}(z)$ is a section of the tautological line bundle on the complex $13$-ball\n$B(L) = \\lbrace v \\in L \\otimes_{\\mathbb{Z}[e^{2 \\pi i\/3}]} \\mathbb{C} \\colon v^2 < 0 \\rbrace\/\\mathbb{C}^*$.\nThus $E^{(l)}(z)$ are meromorphic automorphic forms of type $U(13,1)$\ninvariant under $\\op{Aut}(L)$ and having poles along\nthe mirrors of the complex hyperbolic reflection group $R(L)$. So \n\\begin{equation*}\nz \\mapsto [E^{(30)}(z)\/E^{(36)}(z) : E^{(36)}(z)\/ E^{(42)}(z) : E^{(42)}(z)\/E^{(48)}(z) : \\dotsb]\n\\end{equation*}\nis a meromorphic map from the ball quotient \n$ (B(L) - \\lbrace \\text{mirrors of $R(L)$} \\rbrace) \/ \\op{Aut}(L)$ to the projective space.\nAllcock's monstrous proposal states that the orbifold fundamental group of this ball quotient\nsurjects onto the monster; see \\cite{Allcock-monstrous}, \\cite{AB-braidlike} and\nthe references in there for more details. The calculation of Fourier series of $E^{(l)}(z)$ is\ncomplicated by the fact that the group of translations is a discrete Heisenberg group\nrather than an abelian group. We think that the analysis of $E(z,s)$ is a warm-up exercise\nfor studying the functions $E^{(l)}(z)$.\n\\par\n{\\bf Acknowledgement:} I would like to thank Prof. Richard Borcherds\nfor many stimulating conversations at the beginning of this work.\nI would like to thank Prof. Eric Weber and Prof. Manjunath Krishnapur for their help in a \ncouple of proofs.\n\\section{Some exponential sums related to even lattices}\n\\label{section-j}\n\\begin{topic}{\\bf Notation: }\n\\label{topic-notation}\nWe fix some notation that will be used throughout this section.\nLet $p$ be a prime. Let $q = p^r$ for some integer $r \\geq 1$.\nLet $M$ be a nonsingular even lattice of rank $m$.\nExcept for lemma \\ref{l-bound-on-magnitude-of-j},\nwe do not assume in this section that $M$ is necessarily positive definite.\n Let $n, d \\in \\mathbb{Z}$ and $n \\geq 1$.\nWe abbreviate $M\/n = M\/n M$.\nWe define\n\\begin{equation*}\nM_n(d) = \\lbrace l \\in M\/n : l^2\/2 \\equiv d \\bmod n \\rbrace.\n\\end{equation*}\nLet $\\e(x) = \\exp(2 \\pi i x)$. Let\n$\\lambda \\in M$. The first goal of this section is to prove theorem \\ref{th-j} which lets us\ncalculate some exponential sums of the form \n\\begin{equation*}\nj_{\\lambda, n}^M(d) = j_{\\lambda, n}(d) = \\sum_{l \\in M_n(d)} \\e (\\ip{l}{\\lambda}\/n)\n\\end{equation*}\nin terms of Kloosterman sums \n\\begin{equation*}\nS(a,b, n) = \\sum_{r \\in \\mathbb{Z}\/n : \\op{gcd}(r,n)=1 } \\e((ar + b \\bar{r})\/n)\n\\end{equation*}\nand Jordan totient function $J_k(n) = n^k \\prod_{p \\mid n} (1 - p^{-k} )$. \nAs usual, $\\bar{r}$ is the inverse of $r$ modulo $n$, \nthat is, $\\bar{r}$ denotes any integer such that $r \\bar{r} \\equiv 1 \\bmod n$.\nIn the definition of $J_k(n)$, the product is over all prime divisors of $n$.\nLet $j_{\\lambda, n} = j_{\\lambda, n}(1)$.\nAs a consequence of \\ref{th-j}, we obtain bounds on $\\abs{j_{\\lambda, n}}$\nin \\ref{l-bound-on-magnitude-of-j} and\nshow in \\ref{l-analytic-continuation-of-Dirichlet-series-of-j}\nthat the Dirichlet series $\\sum_{n = 1}^{\\infty} j_{\\lambda, n} n^{-s}$\ncan be analytically continued to a meromorphic function on the whole complex\nplane.\n\\par\nGiven $\\lambda \\in M$, let $c(\\lambda)$ be the largest positive integer such that\n$\\lambda \\in c(\\lambda) M$. So $\\lambda$ is a primitive vector of $M$\nif and only if $c(\\lambda) = 1$.\nWe say $n$ divides $\\lambda$ (resp. $n$ and $\\lambda$ are relatively prime)\nif $n$ divides $c(\\lambda)$ (resp. if $n$ and $c(\\lambda)$ are relatively prime).\nIf $p$ is a prime, we let $v_p(\\lambda)$ be the $p$-valuation of $c(\\lambda)$.\n\\end{topic}\n\\begin{theorem} Assume the setup of \\ref{topic-notation}.\n\\par\n(a) If $r, n$ are relatively prime positive integers, then\n$M_{n r}(d) \\simeq M_r(d) \\times M_n(d)$ and \n$j_{\\lambda, nr}(d) = j_{\\bar{n} \\lambda, r}(d) j_{\\bar{r} \\lambda, n}(d)$.\n\\par\n(b) If $n \\mid \\lambda$, then $j_{\\lambda, n} (d) = j_{0, n} (d) = \\abs{M_n(d)}$.\nIf $p \\nmid d$, \nthen $j_{ p \\lambda, p q}(d) = p^{m-1} j_{\\lambda, q} (d)$.\n\\par\n(c) Assume that $M$ is self-dual and that $n$ and $\\lambda$ are relatively prime.\nThen\n$j_{\\lambda, n}(d) = n^{(m\/2) - 1} S(d, \\lambda^2\/2, n)$. \n\\par\n(d) Assume $M$ is self-dual. Then\n$ j_{0, n}(1) = n^{(m\/2) - 1} J_{m\/2}(n)$.\n\\label{th-j}\n\\end{theorem}\nFor later application,\nwe need to calculate $j_{\\lambda,n} = j_{\\lambda, n}(1)$ for an even self-dual\nlattice. Theorem \\ref{th-j} is sufficient for this purpose. \nFirst, because of part (a), it suffices to calculate\n$j_{\\lambda, q}$ when $q$ is a prime power.\nNext, part (b) lets us reduce to the case when $p \\nmid \\lambda$ or $\\lambda = 0$.\nFinally, parts (c) and (d) handle these two cases respectively.\nProofs of parts (c) and (d) require calculation of some quadratic Gauss sums.\nThese calculations have been moved back to appendix \\ref{appendix-gauss-sums}.\n\\begin{proof}[proof of \\ref{th-j}(a)]\nChoose $\\bar{n}, \\bar{r} \\in \\mathbb{Z}$ such that $\\bar{n} n + \\bar{r} r = 1$.\nFrom Chinese remainder theorem we have mutually inverse isomorphisms \n$\\pi: M\/nr \\to M\/n \\times M\/r$ and $\\phi: M\/n \\times M\/r \\to M\/nr$ given by\n\\begin{equation*}\n\\pi( l \\bmod nr) = (l \\bmod n, l \\bmod r), \\;\\;\n\\phi(l_1 \\bmod n , l_2 \\bmod r) = ( \\bar{n} n l_2 + \\bar{r} r l_1) \\bmod n r\n\\end{equation*}\nfor $l , l_1, l_2 \\in M$.\nPick $v = (l_1 \\bmod n, l_2 \\bmod r) \\in M_n(d) \\times M_r(d)$.\nThen $\\phi(v) = (l \\bmod n r)$ where $l = \\bar{n} n l_2 + \\bar{r} r l_1$. We have\n\\begin{equation*}\nl^2\/2\n \\equiv \\bar{n}^2 n^2 (l_2^2\/2) + \\bar{r}^2 r^2(l_1^2\/2) \n \\equiv \\bar{n}^2 n^2 d + \\bar{r}^2 r^2 d \n\\equiv d \\bmod n r.\n\\end{equation*}\nSo $\\phi(v) \\in M_{ n r }(d)$. This proves $\\phi(M_n (d) \\times M_r(d)) \\subseteq M_{n r}(d)$.\nClearly $\\pi (M_{n r}(d)) \\subseteq M_n (d) \\times M_r(d)$. So\nrestrictions of $\\pi$ and $\\phi$ yield mutually inverse\nisomorphisms between $M_{n r}(d)$ and $M_n (d) \\times M_r(d)$. \nNow let \n$(l_1,l_2)$\nrun over $M_n(d) \\times M_r(d)$.\nThen $\\phi(l_1,l_2)$\nruns over $M_{n r}(d)$.\nIt follows that\n\\begin{equation*}\nj_{\\lambda, n r}(d) \n= \\sum_{l_1 , l_2} \\e\\bigl(\\tfrac{\\ip{ \\bar{n} n l_2 + \\bar{r} r l_1 }{\\lambda}}{n r} \\bigr) \n= \\sum_{l_1 , l_2} \\e \\bigl(\\tfrac{\\ip{ l_2}{\\bar{n} \\lambda}}{ r} \\bigr) \n\\e\\bigl( \\tfrac{\\ip{l_1}{ \\bar{r} \\lambda}}{n} \\bigr)\n= j_{\\bar{n} \\lambda, r}(d) j_{\\bar{r} \\lambda, n} (d).\n\\qedhere\n\\end{equation*}\n\\end{proof}\nPart (b) of theorem \\ref{th-j} can be proved using a Hensel's lemma type argument. We isolate\nthis argument in the proof of the next lemma.\n\\begin{lemma} Assume the setup of \\ref{topic-notation}.\nFurther, assume that $M$ is self-dual and $p \\nmid d$. Then\nthe natural projection $\\pi: M\/p q \\to M\/q$ maps\n$M_{p q } (d)$ onto $M_{q} (d)$.\nFor each $\\bar{l} \\in M_q(d)$, there exists $p^{m-1}$ elements \n $l \\in M_{p q}(d)$ such that $\\pi(l) = \\bar{l}$. \n\\label{l-Hensel}\n\\end{lemma}\n\\begin{proof} Clearly $\\pi(M_{p q}(d) ) \\subseteq M_q(d)$.\nGiven $\\bar{u} \\in M_{q} (d)$, \nfix $u \\in M$ such that $\\bar{u} = u \\bmod qM$.\nThe lifts of $\\bar{u}$ to $M\/p q$ are of the form\n$(u + q x) \\bmod p q M$\n where $x$ runs over a full set of coset representatives for $M\/p$.\nSince $\\bar{u} \\in M_{q}(d)$, we have $\\tfrac{1}{2} u^2 \\equiv d \\bmod q$. \nNote that $(u + q x) \\bmod p q M \\in M_{p q}(d)$ if and only if\n\\begin{equation*}\n \\tfrac{1}{2} u^2 + q \\ip{u}{x} \\equiv d \\bmod p q, \n\\end{equation*}\nor equivalently,\n\\begin{equation}\n\\ip{u}{x} \\equiv q^{-1} (d - \\tfrac{1}{2} u^2) \\bmod p.\n\\label{eq-ipux}\n\\end{equation}\nSince $M$ is self-dual, the bilinear form $\\ip{x}{y} \\bmod p$ on the $\\mathbb{F}_p$ -vector space\n$M\/p$ is non-degenerate. \nNote that $(u \\bmod p M)$ is a nonzero vector in $M\/p$ Since\n$\\tfrac{1}{2} u^2 \\equiv d \\bmod q$ and $p \\nmid d$.\nSo the set of all possible $x \\in M\/p$ satisfying\nequation \\eqref{eq-ipux} forms an affine hyperplane in the $m$ dimensional \n$\\mathbb{F}_p$ -vector space $M\/p$. \n\\end{proof}\n\\begin{proof}[proof of theorem \\ref{th-j}(b)]\nThe first claim is obvious from the definition of $j_{\\lambda, n}(d)$. \nThe second claim follows from lemma \\ref{l-Hensel}, since\n\\begin{equation*}\nj_{p \\lambda, p q}(d) \n= \\sum_{ l \\in M_{p q}(d)} \\e( \\ip{l}{\\lambda}\/ q ) \n= p^{m - 1} \\sum_{ \\bar{l} \\in M_{q} (d)} \\e( \\ip{\\bar{l}}{\\lambda}\/ q ) \n= p^{m-1} j_{\\lambda, q} (d) .\n\\end{equation*}\n\\end{proof}\nGiven a function $f : \\mathbb{Z}\/n \\to \\mathbb{C}$, let \n$(\\mathcal{F} f)(x) = \\sum_{y \\in \\mathbb{Z}\/n} f(y) \\e(x y\/n)$\ndenote its Fourier transform.\nLemma \\ref{l-Fj} calculates the Fourier transform of \n$d \\mapsto j_{\\lambda, q}(d)$.\nThen we recover $j_{\\lambda, q}(d)$ by using the Fourier\ninversion formula $(\\mathcal{F}^2 f)(x) = n f(-x)$.\n\\begin{lemma}\nAssume the setup of \\ref{topic-notation} and that $M$ is self-dual.\n Let $\\lambda \\in M$.\n\\par\n(a) We have \n\\begin{equation}\n\\mathcal{F} j_{\\lambda, q}(c) =\n\\sum_{l \\in M \/q } \\e ( (2 \\ip{l}{\\lambda} + c l^2)\/2q ).\n\\label{eq-Fj}\n\\end{equation}\n\\par\n(b) Assume $p \\nmid \\lambda$. If $p \\mid c$, then $\\mathcal{F} j_{\\lambda, q} (c) = 0$. \nIf $p \\nmid c$, then\n\\begin{equation}\n\\mathcal{F} j_{\\lambda, q} (c) = \\e( - \\bar{c} \\lambda^2\/2 q ) q^{m\/2}.\n\\label{eq-Fj-rel-prime-case}\n\\end{equation}\n\\label{l-Fj}\n\\end{lemma}\n\\begin{proof} \n(a) Note that $M\/q$ is the disjoint union of $M_q(1), M_q(2), \\dotsb, M_q(q)$.\nPart (a) now follows since\n\\begin{equation*}\n \\mathcal{F} j_{\\lambda, q}(c) \n= \\sum_{d = 1}^q \\sum_{\\substack{l \\in M \/q \\\\ l^2\/2 \\equiv d \\bmod q}}\n\\e\\bigl( \\tfrac{\\ip{l}{\\lambda}}{q} \\bigr) \\e \\bigl(\\tfrac{c d}{q} \\bigr)\n= \\sum_{l \\in M \/q } \n\\e\\bigl( \\tfrac{\\ip{l}{\\lambda}}{q} + \\tfrac{c l^2}{2q} \\bigr).\n\\end{equation*}\n\\par\n(b) If $\\op{gcd}(c,q) = 1$, we can complete squares in equation \\eqref{eq-Fj} to get\n\\begin{equation*}\n \\mathcal{F} j_{\\lambda, q}(c) \n= \\sum_{l \\in M\/q } \\e\\bigl( \\tfrac{-\\bar{c} \\lambda^2}{2q } + \\tfrac{c(l + \\bar{c}\\lambda)^2}{2q} \\bigr)\n= \\e \\bigl( \\tfrac{-\\bar{c} \\lambda^2}{2q} \\bigr) \\sum_{l \\in M \/q } \\e \\bigl( \\tfrac{c l^2}{2q} \\bigr).\n\\end{equation*}\nEquation \\eqref{eq-Fj-rel-prime-case} now follows from the lemma \\ref{l-theta-q} in the appendix.\n\\par\nNow assume that $c = p c_1$ for some $c_1 \\in \\mathbb{Z}$.\nLet $u$ run over a set of coset representatives for $M\/p^{r-1} $\nand $x$ run over a set of coset representatives of $M\/ p $. Then $(u + p^{r-1} x)$ runs over\nover a set of coset representatives of $M\/p^r$. \nSo, substituting $ l = (u + p^{r-1} x)$ in equation \n\\eqref{eq-Fj}, we \nobtain\n\\begin{align*}\n \\mathcal{F} j_{ \\lambda, p^r}(c) \n = \\sum_{u \\in M\/p^{r-1} } \n\\e \\bigl( \\tfrac{\\ip{u}{\\lambda}}{p^r} + \\tfrac{c_1 u^2}{2p^{r-1} } \\bigr) \\sum_{x \\in M\/p} \n\\e \\bigl( \\tfrac{\\ip{x }{\\lambda}}{ p} \\bigr). \n\\end{align*}\nIf $p \\nmid \\lambda$, then\n$x \\mapsto \\ip{x }{\\lambda}\/ p$ is a nontrivial character on $M\/p$, since $M$ is self-dual. \nSo the inner sum in the last expression vanishes.\n\\end{proof}\n\\begin{proof}[proof of theorem \\ref{th-j}(c)]\nWhen $n = q$ is a prime power, the formula for $j_{\\lambda, n}(d)$ follows from \nlemma \\ref{eq-Fj-rel-prime-case} using\nFourier inversion. The\ngeneral case follows from this using multiplicativity property of $j_{\\lambda, n}$ \nand Kloosterman sums (see \\ref{th-j}(a) and \\cite{IK:ANT} respectively).\n\\end{proof}\n\n\\begin{proof}[proof of theorem \\ref{th-j}(d)]\nBoth sides of the formula we want to prove are multiplicative. So it is enough to \nprove the formula when $n = q = p^r$ is a prime power.\nFrom equation \\eqref{eq-Fj} and lemma \\ref{l-theta-q},\nwe have\n\\begin{equation*}\n\\mathcal{F} j_{0,p}(c) = \\sum_{l \\in M\/p } \\e( cl^2\/ 2p) \n= \\begin{cases} p^{ m\/2} & \\text{\\; if \\;} 1 \\leq c < p \\\\\np^m & \\text{\\; if \\;} c = p.\n\\end{cases}\n\\end{equation*}\nUsing Fourier inversion, we find \n\\begin{equation*}\np j_{0,p} (1) \n= \\sum_{c = 1}^p \\mathcal{F} j_{0,p}(c) \\e(-c\/p) \n= (p^m - p^{m\/2}). \n\\end{equation*}\nLemma \\ref{l-Hensel} implies\n$j_{0,pq}(1) = \\abs{J_{p q}(1)} = p^{m-1} \\abs{J_{q }(1)} = p^{m-1} j_{0,q}(1)$.\nPart (d) of \\ref{th-j} follows from this.\n\\end{proof}\nRecall the notation $j_{\\lambda,n} = j_{\\lambda,n}(1)$.\n\\begin{lemma} \nAssume that $M$ is self-dual.\nLet $u = \\min \\lbrace v_p(n), v_p(\\lambda) \\rbrace > 0$.\nThen \n\\begin{equation*}\n j_{\\lambda, n} = \\begin{cases}\np^{(m-1)u} j_{p^{-u} \\lambda, p^{-u} n} &\\text{\\; if \\;} v_p(\\lambda) < v_p(n), \\\\\n p^{(m-1)u} (1 - p^{-m\/2}) j_{p^{-u} \\lambda, p^{-u} n} & \\text{\\; if \\;} v_p(\\lambda) \\geq v_p(n).\n\\end{cases}\n\\end{equation*}\n\\label{l-cancel-common-factor-in-j}\n\\end{lemma}\n\\begin{proof} This is a straight forward computation using \\ref{th-j} (a), (b) in the first case and \\ref{th-j} (a), (d) in the second case.\n\\end{proof}\n\\begin{lemma} Assume $M$ is positive definite and self-dual.\nFix $\\epsilon > 0$. Then there exists a constant $C_{\\epsilon}$ such \nthat\n\\begin{equation*}\n\\abs{j_{n, \\lambda}} < C_{\\epsilon} \\abs{\\lambda}^{(m-1)\/2 } n^{\\epsilon + (m-1)\/2 }\n\\text{\\; for all \\;} n \\in \\mathbb{N} \\text{\\; and all \\;} \\lambda \\in M - \\lbrace 0 \\rbrace.\n\\end{equation*}\n\\label{l-bound-on-magnitude-of-j}\n\\end{lemma}\n\\begin{proof}\nSince $M$ is positive definite and integral, $\\lambda^2 \\geq \\op{gcd}(n, \\lambda)^2$. So\nit suffices to prove an inequality of the form\n\\begin{equation}\n\\abs{j_{n, \\lambda}} < C_{\\epsilon} (n \\op{gcd}(n, \\lambda))^{(m-1)\/2} n^{\\epsilon}.\n\\label{eq-bound-on-j}\n\\end{equation}\nWe prove \\eqref{eq-bound-on-j} by induction on $\\op{gcd}(n,\\lambda)$.\nIf $\\op{gcd}(n, \\lambda) = 1$, then \\eqref{eq-bound-on-j} follows \nfrom \\ref{th-j}(c) and Weil's bound for Kloosterman sum (see \\cite{IK:ANT}, p. 280)\n\\begin{equation*}\n\\abs{S(a,b,n) } \\leq \\bigl( \\sum\\nolimits_{d \\mid n} 1 \\bigr) \\op{gcd}(a,b,n)^{1\/2} n^{1\/2} \n\\end{equation*}\nand the estimate $\\sum_{d \\mid n} 1 = o( n^{\\epsilon})$ for any $\\epsilon > 0$.\n\\par \nIf $\\op{gcd}(n, \\lambda) > 1$, then choose a prime $p$ dividing $\\op{gcd}(n, \\lambda)$. \nLemma \\ref{l-cancel-common-factor-in-j}\nimplies that \n$\\abs{j_{\\lambda, n} } \\leq p^{(m-1)u} \\abs{j_{p^{-u} \\lambda, p^{-u} n}(d) }$\nwhere $u = \\min \\lbrace v_p(n), v_p(\\lambda)\\rbrace$.\nNow one can induct.\n\\end{proof}\nWe end this section by showing that the Dirichlet series\n\\begin{equation*}\nL(j_{ \\lambda}, s) = \\sum_{n \\in \\mathbb{N}} j_{\\lambda, n} n^{-s}\n\\end{equation*}\ncan be meromorphically continued to $\\mathbb{C}$ when the lattice $M$ is self-dual. \nThis follows from Selberg's theorem on analytic continuation\nof Kloosterman sum zeta function. The specific result we need is\nrecorded in the next lemma. \n\\begin{lemma}\nThe series $\\sum_{n \\colon \\op{gcd}(n, k) = 1} S(a,b, n) n^{-s}$ can be analytically continued to a meromorphic function on the\nwhole $s$-plane. It is analytic for $\\op{Re}(s) > 1$, except possibly for a finite set of poles on the real segment $1 < s < 2$.\n\\label{l-Selberg}\n\\end{lemma}\n\\begin{proof} \n In the notation used in\n equation (3.9) of \\cite{AS:FCM}, the Dirichlet series\n $F(d) = \\sum_{n \\equiv 0 \\bmod d} S(a,b, n) n^{-s}$ is \n equal to $Z(s\/2,a,b, \\chi, \\Gamma_0(N))$ for the trivial multiplier \n system $\\chi = 1$. By results of \\cite{AS:FCM}\n the series $F(d)$ can be analytically continued to a meromorphic function on $\\mathbb{C}$\n (also see \\cite{DI:KS} for details).\nThe lemma follows since\n$\\sum_{n : \\op{gcd}(n,k) = 1 } S(a,b, n) n^{-s} = \\sum_{d \\mid k} \\mu(d) F(d)$.\n\\end{proof}\n\\begin{lemma} \nAssume the setup of \\ref{topic-notation}. Further, assume that $M$ is self-dual.\nLet $T$ be a finite set of primes. Let $\\mathbb{N}_T$ be the set of all positive integers $n$\nsuch that $ p \\nmid n$ for all $p \\in T$.\nLet $\\lambda \\in M - \\lbrace 0 \\rbrace$. Then the Dirichlet series\n\\begin{equation*}\nL_T( j_{\\lambda}, s) = \\sum_{n \\in \\mathbb{N}_T } j_{\\lambda, n} n^{-s}\n\\end{equation*}\ncan be analytically\ncontinued to a meromorphic function on the whole $s$-plane.\nFor $\\lambda = 0$, we have \n$L(j_{ 0} , s) = \\zeta( s + 1 - m)\/ \\zeta( s + 1 - (m\/2))$.\n\\par\nIn particular, $L(j_{\\lambda}, s)$ has analytic continuation as a meromorphic function on \n the whole $s$-plane for all $\\lambda \\in M$.\n\\label{l-analytic-continuation-of-Dirichlet-series-of-j}\n\\end{lemma}\n\\begin{proof}\nLet $P(\\lambda)$ be the set of prime divisors of $\\lambda$. \nWe shall prove the lemma by induction on $\\abs{P(\\lambda) - T }$.\nFirst, assume that $T \\supseteq P(\\lambda)$. \nIf $n \\in \\mathbb{N}_T$, then $n$ and $\\lambda$ are relatively prime.\nSo by theorem \\ref{th-j}(c), we have\n\\begin{equation*}\nL_T( j_{\\lambda}, s) = \\sum_{n \\in \\mathbb{N}_T} S(1, \\lambda^2\/2, n) n^{(m\/2) - 1 - s}.\n\\end{equation*}\nNote that $\\mathbb{N}_T = \\lbrace n \\in \\mathbb{N} \\colon \\op{gcd}(n, k_T) =1 \\rbrace$\nwhere $k_T$ is the product of all the primes in $T$.\nSo \\ref{l-Selberg} proves the lemma in the case $\\abs{ P(\\lambda) - T } = 0$\n\\par\nNow fix a $\\lambda \\in M$ and a finite set of primes $T$ such that\n$P(\\lambda) - T \\neq \\emptyset$. We want to prove the lemma for $(\\lambda , T)$.\nAssume by induction that the lemma holds for all $(\\lambda', T')$\nsuch that $\\abs{ P(\\lambda') - T' } < \\abs{ P(\\lambda) - T} $. \nChoose a prime $p \\in P(\\lambda) - T$. Let $T' = T \\cup \\lbrace p \\rbrace$.\nWrite $\\lambda = p^r \\mu$ with $p \\nmid \\mu$. \nWrite\n$L_T(j_{ \\lambda}, s) = A_{\\infty} + \\sum_{t = 0}^{r} A_t$\nwhere\n\\begin{equation*}\nA_t = \\sum_{n \\in \\mathbb{N}_T : v_p(n) = t} j_{\\lambda, n} n^{-s} \n\\text{\\; and \\;} \nA_{\\infty} = \\sum_{n \\in \\mathbb{N}_T : v_p(n) > r} j_{\\lambda, n} n^{-s}.\n\\end{equation*}\nWe shall argue that each $A_j$ can be analytically continued.\nFist consider the sum $A_t$ with $t \\leq r$. \nSo let $n \\in \\mathbb{N}_T$ such that $v_p(n) = t \\leq r$. Let $ n_1 = p^{-t} n$.\nUsing lemma \\ref{l-cancel-common-factor-in-j}, we obtain\n\\begin{equation*}\nA_t \n= \\sum_{n_1 \\in \\mathbb{N}_{ T' } } j_{0,p^t} \\cdot j_{p^{-t} \\lambda, n_1} (p^t n_1)^{-s}\n= p^{-t s} j_{0,p^t} \\cdot L_{T'} (j_{ p^{-t} \\lambda} , s).\n\\end{equation*}\nNote that $\\abs{ P(p^{-t} \\lambda) - T' }= \\abs{ (P(\\lambda) - T) - \\lbrace p \\rbrace } \n= \\abs{P(\\lambda) - T} - 1$. \nSo $A_t$ has an analytic continuation by induction hypothesis.\n\\par\nNow consider the sum $A_{\\infty}$. Let $n \\in \\mathbb{N}_T$ such that $v_p(n) > r$. \nLemma \\ref{l-cancel-common-factor-in-j} implies\n\\begin{equation*}\nj_{\\lambda, n} \n= p^{(m-1) r} j_{\\mu, p^{-r} n}.\n\\end{equation*}\nWrite $n_2 = p^{-r} n$. Note that as $n$ varies over\n$\\mathbb{N}_T \\cap \\lbrace k \\in \\mathbb{N} \\colon v_{p}(k) > r \\rbrace$,\nthe number $n_2$ varies over $\\mathbb{N}_T - \\mathbb{N}_{T'}$.\nIt follows that\n\\begin{equation*}\nA_{\\infty} \n= p^{(m-1) r} \\sum_{n_2 \\in \\mathbb{N}_{T} - \\mathbb{N}_{T'} } j_{\\mu, n_2} (p^{r} n_2)^{-s}\n= p^{(m - 1 - s) r} ( L_{T} ( j_{ \\mu}, s) - L_{T' } ( j_{ \\mu}, s) ).\n\\end{equation*}\nNote that $\\abs{P(\\mu) - T} = \\abs{P(\\mu) - T' } = \\abs{P(\\lambda) - T} - 1$. So\nby induction hypothesis, it follows that\n$A_{\\infty}$ has analytic continuation. This proves the lemma for $\\lambda \\neq 0$.\n\\par\nFinally, for $\\lambda = 0$, theorem \\ref{th-j}(d) implies\n$L(j_0,s) = L( J_{m\/2}, s + 1 - (m\/2))$.\nSince $\\sum_{d \\mid n} J_{k} (n) = n^{k}$, we have\n$L(J_{k}, w) = \\zeta( w - k)\/\\zeta(w)$. \n\\end{proof}\n\\section{The even self-dual lattice of signature (25,1)}\n\\label{section-II251}\nIn this section, we collect a few facts about the Leech lattice and $\\mathrm{II}_{25,1}$ that\nwe would need for our calculations in section \\ref{section-Fourier-expansion}. \nFor more details see \\cite{C:Aut26} or \\cite{REB:Thesis}.\n Our sign conventions about Lorentzian \nlattices are consistent with \\cite{REB:Thesis}.\n\\begin{topic}{\\bf Leech cusp and Weyl chamber ``containing\" a Leech cusp:} \n\\label{t-C}\nAs in the introduction, Let $L = \\Lambda \\oplus \\mathrm{II}_{1,1} \\simeq \\mathrm{II}_{25,1}$\nwhere $\\Lambda$ is the Leech lattice and $\\mathrm{II}_{1,1}$ is a hyperbolic cell.\nWrite $V = L \\otimes \\mathbb{R}$. \nIf $l \\in L$, we shall write $l = (\\lambda; m,n)$ with $\\lambda \\in \\Lambda$ and $m,n \\in \\mathbb{Z}$.\nThen $l^2 = \\lambda^2 - 2 m n$. \nLet $\\rho = (0; 0, 1) \\in L$. Then $\\rho$ is a primitive norm zero vector in $L$ such that \n$\\rho^{\\bot}\/\\rho \\simeq \\Lambda$, that is, $\\rho$ is a Leech cusp. \nIf $v \\in V$, we define its {\\it height} by\n\\begin{equation*}\n\\op{ht}(v) = - \\ip{v}{\\rho}.\n\\end{equation*}\nSo the height of $( *; h, *)$ is equal to $h$.\nThe automorphism group of $L$ acts on the hyperbolic space $\\mathcal{H}^{25}$.\nAs a concrete model of $\\mathcal{H}^{25}$, we take\n\\begin{equation*}\n\\mathcal{H}^{25} = \\lbrace v \\in V \\colon v^2 = -1, \\op{ht}(v) > 0 \\rbrace.\n\\end{equation*}\nSo if $v \\in \\mathcal{H}^{25}$, then it has the form $ v = (x; m, (x^2 + 1)\/2m) $ for some $x \\in \\Lambda \\otimes \\mathbb{R}$ and some $m > 0 $.\nIf $c \\in \\mathbb{R}$, define\n \\begin{equation*}\n B_{\\rho}(c) = \\lbrace v \\in \\mathcal{H}^{25} \\colon \\op{ht}(v) < c \\rbrace.\n \\end{equation*}\nThe sets of the form $B_{\\rho}(c)$ are called {\\it (open) horoballs around $\\rho$}.\nGiven $x, y \\in \\mathcal{H}^{25}$, we say that $x$ is closer to $\\rho$ than $y$ if\n$\\op{ht}(x) < \\op{ht}(y)$. This is equivalent to saying that $d(x, B) < d(y,B)$ where\n $B$ is a small horoball around $\\rho$ that does not meet $\\lbrace x, y \\rbrace$.\n\\par\nA {\\it root} of a lattice is a norm $2$ lattice vector. Let $r$ be a root of $L$.\nWe say that $r$ is a {\\it positive root} if $\\op{ht}(r) > 0$. \nThe hyperplane $r^{\\bot}$ determines a \nhyperplane in $\\mathcal{H}^{25}$, called the {\\it mirror} of the root $r$.\nSince the Leech lattice has no roots, \n$\\rho$ does not lie on the mirror of any root of $L$.\nThe positive roots of $L$ corresponding to the mirrors closest to $\\rho$ are the roots of height $1$. These\nare called the {\\it Leech roots} (or {\\it simple roots}) and are parametrized by the vectors of the Leech lattice.\nFor each $\\lambda \\in \\Lambda$, we have a Leech root \n\\begin{equation*}\ns_{\\lambda} = ( \\lambda; 1, \\tfrac{\\lambda^2}{2} - 1).\n\\end{equation*}\nThere is a unique Weyl chamber $C$ of the reflection group $R(L)$ in $\\mathcal{H}^{25}$ \ncontaining $\\rho$ in its closure.\nConway proved that\n\\begin{equation*}\nC = \\lbrace x \\in \\mathcal{H}^{25} \\colon \\ip{x}{s_{\\lambda}} < 0 \\text{\\; for all \\;} \\lambda \\in \\Lambda \\rbrace.\n\\end{equation*}\nAll these inequalities are necessary to define $C$, that is, \nthe walls of $C$ are in bijection with the mirrors\northogonal to the Leech roots. \n\\end{topic}\n\\begin{lemma}\nThe horoball $B_{\\rho}(1\/\\sqrt{2})$ is contained in the Weyl chamber $C$.\n\\label{l-small-horoballs-fit-in-Weyl-chamber}\n\\end{lemma}\n\\begin{proof}\nTake $v = (\\alpha; h, \\tfrac{\\alpha^2 +1}{2 h} ) \\in B_{\\rho}(1\/\\sqrt{2})$. Then $h^2 < \\tfrac{1}{2}$.\nWriting out $\\ip{s_{\\lambda}}{v}$ and completing square one obtains\n\\begin{equation*}\n\\ip{s_{\\lambda}}{v} = -\\tfrac{h}{2} ( \\lambda - \\tfrac{\\alpha}{h} )^2 + h - \\tfrac{1}{2 h}.\n\\end{equation*}\nNow $h^2< \\tfrac{1}{2}$ implies $h - \\tfrac{1}{2h} < 0$. So\n$\\ip{s_{\\lambda}}{v} < 0$ for all $\\lambda \\in \\Lambda$.\n\\end{proof}\n\\begin{topic}{\\bf The translations:}\nLet $\\mathbb{T}_{\\mathbb{R}}$ be the automorphisms of $V$ that fix $\\rho$ and act trivially on\n$\\rho^{\\bot}\/\\rho$.\nOne verifies that $\\mathbb{T}_{\\mathbb{R}} = \\lbrace T_v : v \\in \\Lambda \\otimes \\mathbb{R} \\rbrace$ where \n\\begin{equation*}\nT_v(l; a,b) = (l + av; a, b + \\ip{l}{v} + av^2\/2).\n\\end{equation*}\nThe group $\\mathbb{T}_{\\mathbb{R}}$ is naturally isomorphic to the additive group\n$\\Lambda \\otimes \\mathbb{R}$ via $ T_v \\mapsto v$.\nLet \n\\begin{equation*}\n\\mathbb{T} = \\mathbb{T}_{\\mathbb{R}} \\cap \\Aut(L) = \\lbrace T_{\\mu} : \\mu \\in \\Lambda \\rbrace,\n\\end{equation*}\n that is,\n$\\mathbb{T}$ consists of the automorphisms of $L$ that fix $\\rho$ and act trivially on $\\rho^{\\bot}\/\\rho$. \n\\end{topic}\n\\begin{topic}{\\bf The orbits of the continuous group of translations acting on $V$: }\n\\label{topic-continuous-group-of-translations}\nThe orbits of action of $\\mathbb{T}_{\\mathbb{R}}$ on $V$ are as follows:\nThe subspace $\\rho^{\\bot}$ splits into the zero dimensional orbit $\\lbrace \\rho \\rbrace$\nand one dimensional orbits $(l,0,0) + \\mathbb{R} \\rho$, for\n$l \\in \\Lambda - \\lbrace 0 \\rbrace$.\nOn the complement of $\\rho^{\\bot}$, the action of $\\mathbb{T}_{\\mathbb{R}}$ is free, so we have\n$24$ dimensional orbits. \nFor each $h, k \\in \\mathbb{R}$, we have a free orbit $S_{k,h}$ defined by\n\\begin{equation*}\nS_{k,h} = \\mathbb{T}_{\\mathbb{R}}( 0;h,k\/2h)= \\lbrace v \\in V: v^2 = -k , \\op{ht}(v) = h \\rbrace.\n\\end{equation*}\nThe map \n\\begin{equation}\n\\phi: \\Lambda \\otimes \\mathbb{R} \\to S_{k,h} \\text{\\; defined by \\;} \\phi(v) = T_v(0;h,k\/2h)\n\\label{eq-phi}\n\\end{equation}\nis an isomorphism with inverse given by $ (l, h, *) \\mapsto l\/h$. \nUnder this identification, the action of $\\mathbb{T}$ on $S_{ k , h}$ corresponds to \nthe action of $\\Lambda$ on $\\Lambda \\otimes \\mathbb{R}$ by translation.\nSo $S_{k,h }\/\\mathbb{T}$ can be identified with the torus $(\\Lambda \\otimes \\mathbb{R})\/ \\Lambda$.\n\\end{topic}\n\\begin{topic}{\\bf The orbits of the discrete group of translations acting on the roots: }\n\\label{topic-discrete-group-of-translations}\nSince $\\Lambda$ does not have any roots, all roots of $L$ have nonzero height.\nFix $n \\geq 1$.\nA root $r \\in L(2)$ of height $n$ has the form $r = (l; n, \\tfrac{l^2\/2 - 1}{n} )$.\nSo, given $l \\in \\Lambda$, there is a root of the form $(l; n , *)$\nif and only if $l^2\/2 \\equiv 1 \\bmod n$. \nAs in \\ref{topic-notation}, we write\n\\begin{equation*}\n\\Lambda_n(1) = \\lbrace \\bar{l} \\in \\Lambda\/n : \\bar{l}^2\/2 \\equiv 1 \\bmod n \\rbrace.\n\\end{equation*}\nFor each coset $\\bar{l} \\in \\Lambda_n(1)$,\nlift $\\bar{l}$ to a lattice vector $l$ and fix once and for all a root\n$r_{n,\\bar{l}} =(l;n,*)$ (recall: the last coordinate of $r_{n , \\bar{l}}$ is determined by\nthe condition $r_{n,\\bar{l}}^2 = 2$).\nNow one verifies easily that the set of roots of height $n$\nis the disjoint union of the free $\\mathbb{T}$-orbits $\\mathbb{T} r_{n, \\bar{l}}$,\nparametrized by $\\bar{l} \\in \\Lambda_n(1)$. \n\\end{topic}\nThe lemma below is immediate corollary of \\ref{l-small-horoballs-fit-in-Weyl-chamber}.\n\\begin{corollary}\nLet $h, k \\in \\mathbb{R}_{>0}$ such that $ 2 h^2\/k < 1$. If $z \\in S_{k,h}$, then \n$z\/\\sqrt{k}$ is contained in the Weyl chamber $C$.\n\\label{l-choice-of-h-k}\n\\end{corollary}\n\\section{Fourier series expansion of the Poincare series}\n\\label{section-Fourier-expansion}\n\\begin{topic}\n\\label{topic-setup-for-Fourier-series-computation}\nWe continue with the setup of section \\ref{section-II251}. \nFix a branch of logarithm on positive\nhalf plane so that $(-1)^s = \\e(s\/2)$.\nAs stated in the introduction,\nthe infinite series $E(z,s) = \\sum_{ r \\in L(2)} \\ip{r}{z}^{-s}$ \nconverges for $\\op{Re}(s) > 25$ and defines an analytic function in $s$\ninvariant under $\\Aut(L)$. \nFix, $k, h > 0$ such that \n\\begin{equation*}\n\\nu = 2h^2\/k < 1.\n\\end{equation*}\nThen \\ref{l-choice-of-h-k} implies that $S_{k,h}$ does not meet any hyperplane\northogonal to the roots of $L$.\nSo the restriction of $E(z,s)$ to\n$S_{k,h} \\simeq \\Lambda \\otimes \\mathbb{R} \\simeq \\mathbb{R}^{24}$ (see \\ref{topic-continuous-group-of-translations})\nis well defined and invariant under the group of translation $\\mathbb{T} \\simeq \\Lambda \\simeq \\mathbb{Z}^{24}$.\nThus $E(z,s)$ defines a function on the torus $S_{k ,h}\/\\mathbb{T}$.\nWe shall now calculate the Fourier series of this function.\nWe call this the Fourier series of $E(z,s)$ at the Leech cusp $\\rho$.\nThe Fourier series has the form:\n\\begin{equation}\nE( z ,s) = \\sum_{\\lambda \\in \\Lambda} a_{\\lambda} (k,h,s)\n\\e ( \\ip{ \\lambda}{ \\phi^{-1}(z)} ) \n\\label{eq-Fourier-series-of-E}\n\\end{equation}\nwhere $\\phi : \\Lambda \\otimes \\mathbb{R} \\xrightarrow[]{\\sim} S_{k,h}$ is as in equation \\eqref{eq-phi}.\nThe Fourier coefficients are given by\n\\begin{equation*}\na_{\\lambda}(k, h ,s) \n= \\int_{v \\in (\\Lambda \\otimes \\mathbb{R})\/\\Lambda} E(\\phi(v),s) \n\\e ( -\\ip{\\lambda}{v} ) d v.\n\\end{equation*}\n\\end{topic}\nTheorem \\ref{th-fourier-coeffients} gives a formula for the Fourier coefficients. The formula\ninvolves the character sums $j_{n, \\lambda} = j_{n, \\lambda}^{\\Lambda}(1)$\nfrom section \\ref{section-j} and the numbers\n\\begin{equation*}\nc_n = \\sqrt{\\tfrac{k}{h^2} - \\tfrac{2}{n^2} }.\n\\end{equation*}\nNote that our choice $2 h^2\/k < 1$ ensures that $c_n^2 > 0$ for all $n \\geq 1$.\n\\begin{theorem}\nWrite $a^*_{\\lambda}(k, h, s) = (h\/2 \\pi)^s \\Gamma(s) a_{\\lambda}(k,h,s)$.\nThen one has\n\\begin{equation*}\n a_{\\lambda}^*(k,h,s) \n=\n\\begin{cases} 2 (1 + \\e(-\\tfrac{s}{2} ))\\abs{\\lambda}^{s - 12}\n\\sum_{n \\geq 1} j_{\\lambda, n} n^{-s} c_{n}^{12-s} K_{12 - s}(2 \\pi \\abs{\\lambda} c_{n} ) \n& \\text{ if \\;} \\lambda \\neq 0, \\\\\n& \\\\\n(1 + \\e(-\\tfrac{s}{2} )) \\Gamma(s - 12) \\pi^{12 - s} \\sum_{n \\geq 1} j_{0, n} n^{-s}\nc_n^{24 - 2 s}\n& \\text{ if \\;} \\lambda = 0.\n\\end{cases}\n\\end{equation*}\n\\label{th-fourier-coeffients}\n\\end{theorem}\n\\begin{proof}\nTake $z \\in S_{k,h}$. Then \\ref{l-choice-of-h-k} implies that $z\/ \\sqrt{k}$ belongs to the Weyl chamber $C$.\nSo by Conway's theorem (see \\cite{REB:L}) $z$ and $\\rho$ are on the same side of $r^{\\bot}$ for every root $r \\in L(2)$.\nSo if $r$ is a positive root, then $-\\ip{r}{z} > 0$.\nUsing this and the decomposition of $L(2)$ into $\\mathbb{T}$-orbits from \\ref{topic-discrete-group-of-translations}, we obtain\n\\begin{align*}\nE(z,s)\n= (1 + (-1)^{-s}) \\sum_{r : \\op{ht}(r) > 0} (-\\ip{r}{z})^{-s} \n= (1 + \\e(-\\tfrac{s}{2} )) \n\\sum_{\\substack{ n \\geq 1 \\\\ \\bar{l} \\in \\Lambda_n(1) } }\n\\sum_{T \\in \\mathbb{T}} (-\\ip{T r_{n,\\bar{l}}}{ z})^{-s}.\n\\end{align*}\nSo\n\\begin{equation*}\na_{\\lambda}(k, h ,s) \n=(1 + \\e(-\\tfrac{s}{2} )) \\int_{v \\in (\\Lambda \\otimes \\mathbb{R})\/ \\Lambda} \n \\sum_{n , \\bar{l} } \\sum_{ \\mu \\in \\Lambda}\n(-\\ip{r_{n,\\bar{l}}}{ T_{\\mu} \\phi( v ) })^{-s} \\e( - \\ip{\\lambda}{v + \\mu} ) d v. \n\\end{equation*}\nSince $T_{\\mu} \\phi(v) = \\phi(v + \\mu)$,\nwe can change the above expression into an integral over the whole vector space\n$\\Lambda \\otimes \\mathbb{R}$ to get:\n\\begin{equation*}\na_{\\lambda}(k,h,s)= (1 + \\e(-\\tfrac{s}{2} )) \\sum_{n \\geq 1} \\sum_{\\bar{l} \\in \\Lambda_n(1)}\n\\int_{v \\in \\Lambda \\otimes \\mathbb{R} } (- \\ip{r_{n,\\bar{l}} }{ \\phi(v) })^{-s}\n\\e (- \\ip{\\lambda}{v} ) d v. \n\\end{equation*}\nRecall that $r_{n,\\bar{l}}=(l;n,\\tfrac{l^2 - 2}{2n})$\nand $\\phi(v) = (vh ;h, \\tfrac{ (vh)^2 + k}{2h} )$.\nExpanding the expression for $-\\ip{r_{n,\\bar{l}} }{ \\phi(v) }$\nand completing square we get\n\\begin{equation*}\n-\\ip{r_{n,\\bar{l}} }{ \\phi(v) } = \\tfrac{n h}{2}( c_n^2 +(v -\\tfrac{ l}{n} )^2 )\n\\end{equation*}\nwhere\n$c_n^2 = (\\tfrac{k}{h^2} - \\tfrac{2}{n^2} )$.\nWe substitute $ u = v - \\tfrac{l}{n} $ in the integral to get\n\\begin{equation*}\n a_{\\lambda}(k,h,s)\n= (1 + \\e(-\\tfrac{s}{2} ))\\sum_{n \\geq 1} \\sum_{\\bar{l} \\in \\Lambda_n(1)}\n\\int_{\\Lambda \\otimes \\mathbb{R}} (\\tfrac{n h}{2})^{-s} (c_n^2 + u^2)^{-s}\n\\e (- \\ip{\\lambda}{ u + \\tfrac{l}{n}} ) du \n\\end{equation*}\nor\n\\begin{equation}\na_{\\lambda} (k, h, s)\n=(1 + \\e(-\\tfrac{s}{2} ))(\\tfrac{2}{h})^s \\sum_{n \\geq 1} j_{\\lambda, n} n^{-s}\n\\int_{\\Lambda \\otimes \\mathbb{R}} (c_n^2 + u^2)^{-s}\n\\e(- \\ip{\\lambda}{u}) du\n\\label{eq-integral-expression-for-fourier-coeff}\n\\end{equation}\nwhere\n\\begin{equation*}\nj_{\\lambda, n} = j_{\\lambda, n}^{\\Lambda}\n= \\sum_{l \\in \\Lambda_n(1)} \\e( -\\ip{l}{\\lambda}\/n)\n= \\sum_{l \\in \\Lambda_n(1)} \\e(\\ip{l}{\\lambda}\/n).\n\\end{equation*}\nNote that the two sums are equal since $l \\mapsto -l$ is an involution of $\\Lambda_n(1)$.\n\\par\nAssume $\\lambda \\neq 0$.\nLet $\\op{Gram}(\\Lambda)$ be a Gram matrix of the $\\Lambda$. So\n$u^2 = u' \\op{Gram}(\\Lambda) u$.\nLet $\\mu = \\op{Gram}(\\Lambda)^{1\/2} \\lambda$.\nSo $ \\mu^2 = \\lambda^2$.\nUse the substitution $v = \\op{Gram}(\\Lambda)^{1\/2} u $\nin the integral in \\eqref{eq-integral-expression-for-fourier-coeff} to find\n\\begin{equation*}\nI = \\int_{\\Lambda \\otimes \\mathbb{R}} ( c_n^2 + u^2)^{-s}\n\\e (- \\ip{ u }{\\lambda}) du\n= \\int_{\\mathbb{R}^{24}} (c_n^2 + v^2)^{-s}\n\\e ( -(v \\cdot \\mu) ) dv.\n\\end{equation*}\nThus $I$ is the Fourier transform of a radial function.\nLet $r = \\abs{v}$, $\\xi = v\/\\abs{v}$ and $d v = r^{23} d r d \\omega(\\xi)$, where\n$\\omega$ is the standard measure on $S^{23}$. Changing to polar co-ordinates we have\n\\begin{equation*}\nI = \\int_{r \\in \\mathbb{R}_{+}} r^{23} (c_n^2 + r^2)^{-s} \\int_{\\xi \\in S^{23}} \n\\e( (-r\\abs{\\mu}) (\\xi \\cdot \\mu\/\\abs{\\mu} ) ) d \\omega(\\xi) d r \n\\end{equation*}\nThe integral over the sphere gives a Bessel function of the first kind.\nUsing lemma 9.10.2 of \\cite{AAR:SF}, we evaluate the integral over the sphere to get\n\\begin{equation*}\nI = \\int_0^{\\infty} (c_n^2 + r^2)^{-s} r^{23} 2\\pi J_{11}(2 \\pi \\abs{\\lambda} r )(r\\abs{\\lambda})^{-11} d r. \n\\end{equation*}\nSubstitute $x = 2 \\pi \\abs{\\lambda} r$ to get\n\\begin{equation*}\nI = (2 \\pi \\abs{\\lambda})^{2 s - 12} \\abs{\\lambda}^{-12} \n\\int_0^{\\infty} x^{12} ((2 \\pi \\abs{\\lambda} c_n)^2 + x^2)^{-s} J_{11}(x ) d x.\n\\end{equation*}\nThe last integral can be expressed in terms of the modified Bessel function $K_{\\nu}$. For\n $\\re(s) > 23\/4$, one has\n\\begin{equation*}\nI = \n(2 \\pi \\abs{\\lambda})^{2 s - 12} \\abs{\\lambda}^{-12} \n(2 \\pi \\abs{\\lambda} c_n)^{12 - s} K_{12 - s} (2 \\pi \\abs{\\lambda} c_n) \/ (2^{s-1} \\Gamma(s)).\n\\end{equation*}\n(see the table of Hankel transforms in \n\\cite{EMOT:TOIT}, vol II. p. 24, \nand substitute $y=1$, $\\nu = 11$ and $s = \\mu + 1 $ in formula (20)). \nSubstituting in equation \\eqref{eq-integral-expression-for-fourier-coeff} and\nsimplifying, we get the formula for the Fourier coefficients \nfor $\\lambda \\neq 0$.\n\\par\nFor $\\lambda = 0$, from equation \\eqref{eq-integral-expression-for-fourier-coeff}\nwe get\n\\begin{equation*}\na_0(k, h, s)\n=(1 + \\e(-\\tfrac{s}{2} ))(\\tfrac{2}{h})^s \\sum_{n \\geq 1} j_{0, n} n^{-s}\n\\int_{\\Lambda \\otimes \\mathbb{R}} (c_n^2 + u^2)^{-s} du \\\\\n\\end{equation*}\nBy changing to polar coordinates and changing variable $r = c_n x$, we obtain\n\\begin{equation*}\na_{0}(k,h,s) = (1 + \\e(-\\tfrac{s}{2} ))(\\tfrac{2}{h})^s \\op{vol}(S^{23}) \\sum_{n \\geq 1} j_{0, n} n^{-s} c_n^{24 - 2s} \n\\int_0^{\\infty} x^{23} (1 + x^2)^{-s} dx\n\\end{equation*}\nThe last integral is equal to $ \\Gamma(12) \\Gamma(s -12)\/ 2 \\Gamma(s)$\nfor $\\op{Re}(s) > 12$.\n(see the table of Mellin transforms in \n\\cite{EMOT:TOIT} Vol I, p. 311 and substitute $h = 1, \\alpha = 1$ $s = 24$ in formula (30)).\nSubstituting the value of the integral and $\\op{vol}(S^{23}) = 2 \\pi^{12}\/\\Gamma(12)$\nand simplifying, we get the formula for $a_0( k,h,s)$.\n\\end{proof}\nUsing theorem \\ref{th-fourier-coeffients}, we now prove that\n$E(z,s)$ can be analytically continued. \nThe argument is similar to an argument used to prove analytic continuation\nof real analytic Eisenstein series (see \\cite{B:A}, p. 68-69).\nRecall that we write $\\nu = 2 h^2\/k < 1$.\nNote that $c_n^2 = (2\/\\nu) ( 1 - \\nu\/n^2)$.\n\\begin{theorem}\nThe series $E(z,s)$ can be analytically continued to a meromorphic function\non the half plane $\\op{Re}(s) > 25\/2$. The only poles of $E(z,s)$\nin this half plane comes from the Fourier coefficient $a_0(k,h,s)$.\n\\label{th-analytic-continuation-of-E}\n\\end{theorem}\n\\begin{proof}\nFrom \\ref{l-analytic-continuation-of-Dirichlet-series-of-j}, \nwe know $\\sum_n j_{0, n} n^{-s} = \\zeta(s - 23)\/\\zeta(s - 11)$; so this Dirichlet series\nhas meromorphic continuation to $\\mathbb{C}$. \nLemma \n\\ref{l-analytic-continuation-of-perturbation} \n(applied with $f(n) = j_{0,n}$ and $g(s,y) = (1 - \\nu y^2)^{12 - s}$)\nimplies that\nthe infinite series for $a_0(k,h,s)$ given in theorem \\ref{th-fourier-coeffients} has an \nmeromorphic continuation to $\\mathbb{C}$. Now let $\\lambda \\neq 0$.\nThe bound on $j_{\\lambda,n}$ from lemma \\ref{l-bound-on-magnitude-of-j}\nshows that the infinite series for $a_{\\lambda}(k,h,s)$ in \\ref{th-fourier-coeffients}\nconverges absolutely and uniformly on compacta, for $\\op{Re}(s) >25\/2$. \nSo each Fourier coefficient defines an analytic function on $\\op{Re}(s) > 25\/2$.\nLet $s $ vary on a compact subset $A$ of $\\lbrace z : \\op{Re}(z) > 25\/2 \\rbrace$.\nLet $\\kappa = \\sqrt{2(1 - \\nu)\/\\nu} > 0$. \nNote that $c_n \\geq \\kappa$ for all $n$.\nSo \n\\begin{equation*}\n\\abs{ c_n^{12 - s} } = (1\/c_n)^{\\op{Re}(s - 12)} \\leq (1\/\\kappa)^{\\op{Re}(s - 12)}\n\\text{\\; and \\;} \ne^{- \\pi \\abs{ \\lambda} c_n} \\leq e^{- \\pi \\kappa \\abs{\\lambda}}.\n\\end{equation*}\nThe Macdonald Bessel function $K_s(y)$ decays exponentially as $y \\to \\infty$.\nIt is convenient to use the bound $\\abs{K_s(y)} \\leq e^{-y\/2} K_{\\op{Re}(s)} (2)$\nfor $y > 4$ (see \\cite{B:A}, p. 66). \nFix $\\epsilon >0$ such that $\\op{Re}(s) > 2 \\epsilon + 25\/2$ for all $s \\in A$.\nFrom \\ref{l-bound-on-magnitude-of-j},\nwe know that there exists a constant $C_{\\epsilon}$ such that\n$\\abs{j_{\\lambda,n}} \\leq C_{\\epsilon} \\abs{\\lambda}^{23\/2} n^{23\/2 + \\epsilon}$. \nUsing the expression for $a_{\\lambda}(k,h,s)$ from \\ref{th-fourier-coeffients} and \nthe bound for $\\abs{j_{n,\\lambda}}$, we find\n \\begin{align*}\n \\abs{a_{\\lambda}(k,h,s)}\n\n\n\n\n\n \\leq\n C(s) \\abs{\\lambda}^{\\op{Re}(s) - 1\/2} e^{- \\pi \\kappa \\abs{ \\lambda}}\n \\end{align*}\nwhere \n$C(s)$ is a fixed positive real valued\n continuous function on $ A$ \n that does not depend on $\\lambda$.\n Since $C(s)$ and $\\op{Re}(s)$ \n stay bounded on a compact set,\n$\\abs{a_{\\lambda}(k,h,s)}$ decays exponentially as a function of $\\abs{\\lambda}$ as\n $\\abs{\\lambda} \\to \\infty$. It follows that the Fourier series for $E(z,s)$ given in \n equation \\eqref{eq-Fourier-series-of-E} converges\n absolutely and\n uniformly on compact subsets of $\\op{Re}(s) > 25\/2$.\n \\end{proof}\n\\begin{remark}\n\\label{remark-analytic-continuation-of-Fourier-coefficients}\nLemma \\ref{l-analytic-continuation-of-perturbation} applied with\n$f(n) = j_{n, \\lambda}$ and\n\\begin{equation*}\ng(s, y) = \\sqrt{1 - \\nu y^2}^{12-s} K_{12 - s}( 2 \\pi \\abs{\\lambda} \\sqrt{2\/\\nu} \\sqrt{1 - \\nu y^2} )\n\\end{equation*}\nimplies that each Fourier coefficient $a_{\\lambda}(k,h,s)$ has\nanalytic continuation to a meromorphic function on the whole $s$-plane. But we are not able \nto establish how these extended functions decay as $\\abs{\\lambda} \\to \\infty$.\nSo we are unable to prove analytic continuation of $E(z,s)$ beyond\nthe line $\\op{Re}(s) = 25\/2$.\n\\end{remark}\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\nDoping in organic semiconductors is an ubiquitous phenomenon appearing when donor and acceptor molecules are combined together, and crucially determines the electronic, optical, and transport properties of the resulting materials~\\cite{aziz+07am,yim+08am,ping+10jpcl,gao+13jmcc,ping-nehe13prb,duon+13oe,salz+13prl,mend+13acie, yang+14ami, mend+15ncom,salz+16acr,fuze+16jpcl,jaco-moul17am,duva+18jpcc,beye+19cm}.\nTwo main mechanisms have been identified to be responsible for doping in organic semiconductors: Integer charge transfer (ICT) manifests itself when an electron is transferred from the donor to the acceptor, leading to the formation of ion pairs. \nPartial charge transfer occurs upon electronic hybridization of the frontier orbitals of the donor and the acceptor, giving rise to a charge-transfer complex (CTC). \nWhile a general consensus has been established around these two processes, recent studies showing that ICT and CTC can coexist in the same molecular material~\\cite{jaco+18mh,neel+18jpcl} have reignited the debate about doping in organic semiconductors.\nOne of the open questions in this respect concerns the role of the donor conjugation length in determining the electronic structure and the optical properties of these systems.\nA recent study by Mendez \\textit{et al.}~\\cite{mend+15ncom} shows that ICT appears when the strong electron acceptor 2,3,5,6-tetrafluoro-7,7,8,8-tetracyano-quinodimethane (F4TCNQ) is blended with the poly-3-hexylthiophene (P3HT) polymer, (see also Refs.~\\citenum{yim+08am,ping-nehe13prb,duon+13oe,karp+17mm,scho+17afm}), while CTCs are formed when the same acceptor is coupled to the quarterthiophene oligomer.\nThese findings suggest a fundamental relation between the doping mechanism and the length of the donor chain.\nOn the way to tackle this problem in its general complexity, it is relevant to first understand how the conjugation length of the donor oligomers impacts the electronic structure and the optical excitations of complexes exhibiting only partial charge transfer.\nRationalizing this behavior in the presence of one specific doping mechanism is essential to set the stage for further theoretical and experimental studies.\n\nWe perform a first-principles study based on density-functional theory (DFT) and many-body perturbation theory (MBPT), which offer state-of-the-art accuracy and insight into the electronic and optical properties of the investigated systems~\\cite{blas-atta11apl,baum+12jctc,fabe+14ptrsa,nied+15afm}.\nWe focus on CTCs formed by thiophene oligomers (nT) $p$-doped by F4TCNQ.\nWe consider thiophene chains with an even number of rings ranging from 4T, which is the shortest oligomer typically investigated in experiments~\\cite{aziz+07am,ping+10jpcl,mend+15ncom}, up to 10T.\nAll systems are modeled as molecular dimers \\textit{in vacuo}~\\cite{zhu+11cm,gao+13jmcc,mend+13acie,mend+15ncom}, with the donor and the acceptor molecules in the $\\pi$-$\\pi$ stacking arrangement~\\cite{zhu+11cm,gao+13jmcc}.\nIn this way, we are able to systematically rationalize the role of the donor length on the electronic and optical properties of CTCs which differ from each other only by the size of the nT oligomers.\nStarting from the level alignment of the individual components, we relate the band-gap reduction upon increasing donor length to the variation of the ionization energy of the system.\nIn the analysis of the optical properties we examine how the absorption onset and the spectral density are affected by the size of the nT oligomer.\nWe finally discuss the character of the lowest-energy excitations in terms of electron and hole density and demonstrate that their behavior is critically sensitive to the donor conjugation length.\n\n\\section{Theoretical Background and Computational Details}\nEquilibrium geometries of the individual donor and acceptor molecules as well as of the CTCs are computed from DFT with the all-electron package FHI-aims~\\cite{blum+09cpc} adopting tight integration grids, TIER2 basis sets~\\cite{havu+09jcp}, and the Perdew-Burke-Ernzerhof~\\cite{perd+96prl} generalized-gradient approximation for the exchange-correlation functional.\nVan der Waals interactions, that play an essential role in $\\pi$-$\\pi$ stacked systems, are accounted for by means of the Tkatchenko-Scheffler scheme~\\cite{tkat-sche09prl}. \nAtomic positions are relaxed until the Hellmann-Feynman forces are smaller than 10$^{-3}$ eV\/\\AA{}.\nThe acceptor molecule adsorbed in the middle of the donor chain gives rise to the most stable configuration, in agreement with previous first-principles results~\\cite{zhu+11cm}.\nFurther details and discussion about the geometry of the CTCs are reported in the Supporting Information. \n\nThe electronic structure and the optical properties of the considered systems are computed with the MOLGW code~\\cite{brun+16cpc}.\nGaussian-type cc-pVDZ basis sets~\\cite{brun12jcp} are used including the frozen-core approximation. \nThe resolution-of-identity approximation is also employed~\\cite{weig+02}. \nThe $GW$ approximation in the perturbative approach $G_0W_0$ is adopted to calculate the quasi-particle (QP) correction to the single-particle levels starting from the solution of the Kohn-Sham equations.\nThe range-separated hybrid functional CAM-B3LYP~\\cite{yana+04cpl} is chosen to obtain a reliable starting point for the $G_0W_0$ calculations~\\cite{fabe+13jcp,brun+13jctc}.\nOptical absorption spectra and (bound) electron-hole pairs are computed from the solution of the Bethe-Salpeter equation (BSE), the equation of motion for the two-particle polarizability, $L = L_0 + L_0 \\, K \\, L$~\\cite{brun+16cpc}.\nThe so-called BSE kernel $K$ includes the repulsive exchange term given by the bare Coulomb potential $\\bar{v}$ and the statically screened electron-hole Coulomb attraction $W$. \nFor further details about this formalism, see, \\textit{e.g.}, Refs.~\\citenum{fabe+13jcp,cocc-drax15prb,brun+16cpc}.\nWe solve the BSE in the Tamm-Dancoff approximation (TDA), which does not impact the nature of the excitations analyzed below.\nMinor effects on the excitation energies and oscillator strengths in line with previous works~\\cite{palu+09jcp,baum+12jctc1,rang+17jcp} are discussed in the Supporting Information (see Figure S6 and Table S1).\nThe character and spatial distribution of the electron-hole pairs are obtained from the analysis of the hole and electron densities, which for the $\\lambda^{th}$ excitation are defined as~\\cite{cocc+11jpcl,deco+14jpcc}\n\\begin{equation}\n\\rho_{h}^{\\lambda} (\\textbf{r})=\\sum_{\\alpha \\beta} A_{\\alpha \\beta}^{\\lambda} \\rvert \\phi_{\\alpha}(\\textbf{r})\\lvert^{2}\n\\label{eq:h}\n\\end{equation}\nand\n\\begin{equation}\n\\rho_{e}^{\\lambda} (\\textbf{r})=\\sum_{\\alpha \\beta} A_{\\alpha \\beta}^{\\lambda} \\rvert \\phi_{\\beta}(\\textbf{r})\\lvert^{2},\n\\label{eq:e}\n\\end{equation}\nrespectively.\nThe coefficients $A_{\\alpha \\beta}^{\\lambda}$ are the square of the normalized BSE eigenvectors that act as weighting coefficients of each transition between occupied $(\\phi_{\\alpha})$ and unoccupied $(\\phi_{\\beta})$ QP states contributing to the $\\lambda^{th}$ excitation.\n\n\\section{Results and Discussion}\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.48\\textwidth]{level_alignment_DFT.png}\n\\caption{Energy level alignment between the isolated 4T (left) and F4TCNQ (right) computed from DFT (CAM-B3LYP functional). The HOMO and the LUMO of the resulting complex are shown in the the middle panel. For comparison, the HOMO of the donor and the LUMO of the acceptor are also depicted.}\n\\label{level_alignment}\n\\end{figure}\n\nWe start our analysis by considering the electronic structure of the CTC formed by the 4T oligomer doped by F4TCNQ, depicted in Figure~\\ref{level_alignment}.\nThe energy values defined therein are computed with respect to the vacuum level set to zero.\nBy first inspecting the electronic energy levels of the individual molecules, we notice that they form a \\textit{staggered} or \\textit{type-II} level alignment, with the lowest-unoccupied molecular orbital (LUMO) of the acceptor lying within the gap of the donor.\nIn the resulting 4T-F4TCNQ complex, the highest-occupied molecular orbital (HOMO) is energetically lower than the HOMO of the donor, while the LUMO is higher than the LUMO of the acceptor (see Figure~\\ref{level_alignment}, middle panel).\nIn agreement with the known behavior of CTCs~\\cite{zhu+11cm,mend+13acie,gao+13jmcc,mend+15ncom}, the band gap of the complex is almost 1 eV larger than the energy difference between the HOMO of 4T and the LUMO of F4TCNQ.\nThe HOMO and LUMO of the CTC exhibit the typical features of bonding and anti-bonding states, respectively, due to the $\\pi$-$\\pi$ coupling between the components: The HOMO is seamlessly extended across the donor and the acceptor molecules, while the LUMO distribution is characterized by an evident nodal plane between the 4T and the F4TCNQ (see Figure~\\ref{level_alignment}, middle panel).\nThe character of the frontier orbitals is an intrinsic property of the CTCs, independent of modeling them as dimers \\textit{in vacuo}.\nBonding and anti-bonding states are formed also when the donor\/acceptor complexes are represented by periodic stacks -- see Figure S5 in the Supporting Information, and Ref.~\\citenum{beye+19cm}.\n\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.48\\textwidth]{level_alignment_g0w0.png}\n\\caption{Energy level alignment computed from $G_{0}W_{0}$ on top of DFT (CAM-B3LYP functional) for the 4T-, 6T-, 8T-, and 10T-F4TCNQ complexes and their individual components.}\n\\label{Elec_all}\n\\end{figure}\n\nLevel alignment and band gaps analyzed above are computed from DFT, employing the hybrid functional CAM-B3LYP.\nAn improved picture can be obtained from MBPT including the QP correction computed from $G_0W_0$.\nIn Figure~\\ref{Elec_all} we report the results of these calculations for all the investigated systems, showing that CTCs are formed also upon increasing donor length, with the qualitative trend for the gaps and the frontier levels being unchanged compared to the case of 4T-F4TCNQ discussed above. \nAs expected upon the inclusion of the electronic self-energy, the size of the band gap of the individual molecules and of the resulting CTCs is increased by about 1 eV. \nWe first consider again the case of the 4T-F4TCNQ complex, shown on the top panel of Figure~\\ref{Elec_all}.\nIn the $G_0W_0$ energy levels of both donor and acceptor separately, the LUMO is shifted up by almost 1 eV while the QP correction to the HOMO is one order of magnitude smaller and moves towards more negative energies.\nAs a result, the energy difference between the HOMO of the donor and the LUMO of the acceptor increases from 1.68 eV in DFT (Figure~\\ref{level_alignment}) to 2.87 eV in $G_0W_0$ (Figure~\\ref{Elec_all}).\nThe CTC exhibits a similar although less pronounced behavior, such that the QP gap of 4T-F4TCNQ is equal to 3.56 eV.\nInspecting now Figure~\\ref{Elec_all} in full we notice that the QP gap decreases by about 0.8 eV going from 4T-F4TCNQ to 10T-F4TCNQ.\nThe band gap reduction in the CTCs is less significant than in the isolated thiophene oligomers with increasing length~\\cite{fabi+05jpca,sun-auts14jctc,cocc-drax15prb}, revealing that the hybridization between donor and acceptor cannot be trivially predicted from the level alignment of the individual constituents. \nYet, it is clear from Figure~\\ref{level_alignment} that the decrease of the electronic gaps in the CTCs with varying nT size is mainly due to the up-shift of the HOMO level energy with the LUMO remaining almost constant.\nAs the orbital energies are computed with the vacuum level set to zero, this finding implies that the electron affinity of the CTCs is almost unaffected by the increasing donor length, while the ionization energy is critically influenced by this parameter.\nThis behavior is not surprising, since the examined systems differ only by the length of the donor doped by the same electron acceptor.\nOn the other hand, this result is not obvious \\textit{a priori} considering the hybridized character of the HOMO and the LUMO in all CTCs.\nFrom this analysis it is possible to deduce the first relevant result of this work: The increasing donor conjugation length leads to a sizable decrease in the band gaps of the CTCs (0.8 eV going from 4T-F4TCNQ to 10T-F4TCNQ) which stems almost entirely from the reduction of the ionization energy.\nThe physical rationale of this behavior can be deduced from the spatial distribution of the frontier orbitals (see Figure S4 in the Supporting Information). \nWhile the bonding and anti-bonding character of the HOMO and the LUMO is shared by all CTCs, the delocalization of the highest-occupied orbital is enhanced upon the increasing donor length.\nOn the contrary, even in the longest nT chains, the LUMO remains confined on the F4TCNQ and on the few thiophene rings interacting with the acceptor through $\\pi$-$\\pi$ coupling.\nThe similar trend given by $G_0W_0$ (Figure~\\ref{Elec_all}) and DFT results (Figure~\\ref{level_alignment} and Figure S3 in the Supporting Information) suggests that quantum confinement and electronic hybridization are mainly responsible for the variation of the gap due to the up-shift of the HOMO with almost constant LUMO energies.\nDifferent from dimers formed by the same moiety (see \\textit{e.g.}, Ref.~\\citenum{hoga+13jcp}), here many-body effects, such as the electronic screening and the self-energy, improve the accuracy of orbital energies and electronic gaps, but do not qualitatively alter the picture given by (hybrid) DFT.\n\n\\begin{figure} \n\\centering\n\\includegraphics[width=0.5\\textwidth]{camb3lyp_spectra.png}\n\\caption{Optical absorption spectra of the 4T-, 6T-, 8T-, and 10T-F4TCNQ complexes. The first three bright peaks are labeled as P$_{1}$, P$_{2}$, and P$_{3}$, and the energy of the QP gaps ($E_{gap}$) is marked by vertical dashed lines. }\n\\label{spectra}\n\\end{figure}\n\nWe now turn to the optical properties of the CTCs.\nIn all complexes, the first excitation is dipole-allowed and gives rise to the weak peak labeled P$_1$ (see Figure~\\ref{spectra}).\nRegardless of the donor length, it has very similar energy ranging between 1.1 eV and 1.2 eV and comparable oscillator strength.\nThis behavior can be better understood by analyzing the character of this excitation in terms of orbital contributions.\nIn Table~\\ref{CTC-composition} we report the composition of P$_1$ in the 4T-F4TCNQ and 8T-F4TCNQ complexes (details about the other CTCs are given in the Supporting Information, Table S2) and in Figure~\\ref{h-e} we display the spatial distribution of the corresponding electron and hole densities (Eqs.~\\ref{eq:h} and~\\ref{eq:e}).\nIn all complexes the first excitation stems from the HOMO-LUMO transition, such that the hole and the electron components of P$_1$ exhibit the same character and spatial extension of the frontier orbitals.\nThis has two main consequences: In all systems the lowest-energy excitation has a Frenkel-like character localized within the CTC. \nMoreover, the comparable overlap between the HOMO and the LUMO (and equivalently of the hole and the electron density) in all complexes explains why P$_1$ has almost constant oscillator strength going from 4T-F4TCNQ to 10T-F4TCNQ.\nWe also recall that P$_1$ retains approximately the same energy in all CTCs irrespective of the donor length, in contrast with the decrease of the electronic gaps in the presence of longer nT chains discussed above. \nTo clarify this point, we take advantage of the adopted formalism where optical excitations are computed including the electron-hole correlation.\nIn this framework, the difference between the quasi-particle HOMO-LUMO gap and the excitation energy of P$_1$ can be interpreted as the exciton binding energy ($E_b$).\nBy comparing the excitation energies in Table~\\ref{CTC-composition} (and Table S2 in the Supporting Information) with the QP gaps in Figure~\\ref{Elec_all}, it is evident that $E_b$ decreases as the length of the donor oligomer increases.\nThe binding energy reduction of P$_1$ upon increasing donor length is therefore comparable with the decrease of the band gap (see Figure~\\ref{Elec_all}), such that these two quantities end up compensating each other.\nIn the presence of longer nT chains the Coulomb screening is enhanced and the quantum confinement of the electron-hole pairs is concomitantly reduced.\nThe Coulomb screening is solely embedded in the direct term of the BSE while quantum confinement affects also the bare exchange interaction between the electron and the hole~\\cite{hoga+13jcp,cocc-drax15prb}. \nDue to these intertwined effects and to the reduction of the QP gaps discussed above, P$_1$ appears at the same energy in all the examined CTCs, regardless of the donor conjugation length.\n\n\\begin{table*} \n\\centering\n\\caption{Energies, oscillator strength (OS), and composition in terms of transitions between molecular orbitals of the first three bright excitations in the 4T-F4TCNQ and 8T-F4TCNQ complexes. HOMO and LUMO are abbreviated by H and L, respectively.}\n\\label{CTC-composition}%\n\\begin{tabular}{lcccl}\n\\hline\n& Excitation & Energy [eV]& OS & Composition \\\\ \\hline\n4T-F4TCNQ & P$_1$ & 1.23 & 0.19& H $\\rightarrow$ L \\\\ \n& P$_2$ & 2.02 & 0.19 & H-1 $\\rightarrow$ L \\\\ \n& P$_3$ & 2.72 & 0.10 & H $\\rightarrow$ L+1 \\\\ \n& & & & H-2 $\\rightarrow$ L \\\\ \\hline \n8T-F4TCNQ & P$_1$ & 1.08 & 0.17 & H $\\rightarrow$ L \\\\ \n& P$_2$ & 1.39 & 0.55 & H-1 $\\rightarrow$ L \\\\ \n& P$_3$ & 2.26 & 2.00 & H $\\rightarrow$ L+1 \\\\ \\hline\t\\end{tabular}\n\\end{table*}\n\nAbove the absorption onset a number of peaks appears in all the spectra shown in Figure~\\ref{spectra}.\nIn the following, we focus on the second and third bright excitations, P$_2$ and P$_3$, and analyze their behavior. \nIn the spectrum of 4T-F4TCNQ both P$_2$ and P$_3$ have comparable oscillator strength with respect to each other and to P$_1$. \nThe electron and hole densities of these three excitations are indeed very similar, corresponding in all cases to Frenkel excitons localized within the CTC (see Figure~\\ref{h-e}). \nTheir different composition in terms of single QP transitions (see Table~\\ref{CTC-composition}) reveals that in this complex, where the sizes of the donor and the acceptor are comparable, also the HOMO-1 and the LUMO+1 are hybridized (further details in the Supporting Information, see Figure S4).\nThis picture drawn by our first-principles results is consistent with the experimental data available for the 4T-F4TCNQ complex.\nIn Ref.~\\citenum{mend+15ncom} the optical density of this CTC is reported: The sharp maximum at the absorption onset is followed by a dip and by a broader and more intense peak around 2.4 eV. \nTwo findings are particularly significant: First, the energy separation between the first and second maxima in the measured spectra is consistent with the computed energy difference between P$_1$ and P$_3$ shown in Figure~\\ref{spectra} and Table~\\ref{CTC-composition}. \nSecond, the relative intensity between the two maxima recorded in the experiment follows the same trend resulting in the calculated spectrum. \nAbsolute values of excitation energies and binding energies are not quantitatively comparable between our theoretical results and the experiment, as the CTCs are modeled here by dimers \\textit{in vacuo}. \nThus, only relative trends can be correlated.\nWe also note that in organic films the exciton binding energy may be affected by the crystalline structure of the doped material, which in turns may influence the hole delocalization -- see, \\textit{e.g.}, Ref.~\\citenum{li+17prm}.\n\\begin{figure*} \n\\centering\n\\includegraphics[width=1.0\\textwidth]{hole-elec.png}\n\\caption{Electron (e) and hole (h) densities of the first (P$_{1}$), second (P$_{2}$) and third (P$_{3}$) optically active excitations of each complex. Isosurfaces are plotted with a cutoff of 0.002 \\AA{}$^{-3}$.}\n\\label{h-e}\n\\end{figure*}\n\nThe scenario changes substantially when the CTCs are formed by oligothiophenes with six rings or more: P$_2$ and P$_3$ have significantly higher oscillator strength than P$_1$ (see Figure~\\ref{spectra}) and their character varies accordingly (see Figure~\\ref{h-e}). \nThe orbital composition of these two excitations is analogous to the case of 4T-F4TCNQ (see Table~\\ref{CTC-composition}), demonstrating that the influence of the donor length manifests itself primarily in the character of the orbitals.\nIn these CTCs with enhanced size mismatch between the donor and the acceptor the orbitals below the HOMO and above the LUMO remain localized only on the nT oligomer (see Supporting Information, Figure S4): Being considerably longer than F4TCNQ, 6T, 8T, and 10T have higher density of states around the gap. \nAs a consequence, P$_2$, which stems from the transition between the HOMO-1 and the LUMO in all CTCs, becomes a charge-transfer exciton with the hole localized on the donor and the electron almost completely spread on the acceptor. \nThe enhanced space separation between electron and hole densities is responsible for the relative low oscillator strength of P$_2$ compared to higher-energy peaks.\nOn the other hand, both the electron and the hole in P$_3$ tend to be increasingly localized on the nT chain as the size of the latter grows, giving rise to a Frenkel exciton within the donor. \nThe significant increase of the oscillator strength of P$_3$ upon increasing length of the thiophene oligomers is therefore explained by the larger overlap between the electron and the hole components of this excitation. \nFinally, by inspecting all in all the spectra shown in Figure~\\ref{spectra}, we notice that different from P$_1$, but in line with the behavior of the electronic gaps, both P$_2$ and P$_3$ red-shift upon increasing donor length and so do also the peaks at higher energies.\nThe concomitant rise of their oscillator strength results in increased spectral density in the visible region.\nThe relevance of this result is apparent: Even in case only partial charge transfer occurs in doped organic semiconductors, photo-excited charge carriers can be effectively generated upon efficient absorption of visible light and enhanced electron-hole separation.\n\nBefore concluding, it is instructive to discuss the trends obtained by our results in light of the available experimental findings on nT-F4TCNQ charge-transfer complexes.\nThe optical spectrum measured for 4T-F4TCNQ with 50$\\%$ doping concentration and reported in Ref.~\\citenum{mend+15ncom} is dominated by a bright peak at the onset which is interpreted as a signature of CTC formation.\nThe character of $P_1$ discussed above is fully compatible with this analysis.\nMoreover, the experimental spectrum is characterized by a sharp increase of the oscillator strength about 2 eV above the optical gap, which is also in line with the results shown in Figure~\\ref{spectra}.\nDue to the lack of measurements on CTCs formed by longer oligomers than 4T, it is not possible to extend this comparison to the other systems investigated in this work.\nHowever, considering the tendency of P3HT to bend, forming linear units consisting of at least half a dozen of thiophene rings, it is reasonable to discuss our results in light of measurements performed on P3HT-F4TCNQ blends where partial charge transfer is locally observed~\\cite{jaco+18mh,neel+18jpcl}.\nThe measured spectra reveal a broad absorption band with a maximum between 2.0 eV and 2.5 eV, which is assigned to transitions between hybridized states and thus interpreted as a signature of CTC formation~\\cite{jaco+18mh,neel+18jpcl}. \nThis feature can be related to the peak P$_3$ in the spectra of 8T-F4TCNQ and 10T-F4TCNQ shown in Figure~\\ref{spectra}, taking into account the expected blue-shift of a few hundred meV compared to experiments induced by modelling the CTCs as dimers \\textit{in vacuo}.\nNew experiments on nT-F4TCNQ CTCs with longer donor chains than 4T are certainly needed to corroborate this analysis.\nA quantitative comparison with future measurements will be enabled by the inclusion in the calculated spectra of an effective screening term accounting for the polarizable environment surrounding the doped species in the crystalline phase.\n\n\\section{Summary and Conclusions}\nIn summary, we have investigated from first-principles many-body theory how the donor conjugation length impacts the electronic structure and the optical excitations of CTCs formed by thiophene oligomers doped by the strong electron acceptor F4TCNQ.\nWe have found that all systems are characterized by hybridized frontier orbitals regardless of the size of the thiophene oligomers. \nThe electronic gaps decrease upon increasing donor length solely due to the reduction of the ionization energy, with the electron affinity remaining almost constant in all CTCs. \nThe optical gaps follow a different trend, being at approximately the same energy irrespective of the donor length.\nTogether with the comparable oscillator strength of the first absorption peak, this finding implies that whenever a nT-F4TCNQ charge-transfer complex is formed, it absorbs the same about of light at approximately the same energy independently of the donor length. \nThe size of the nT oligomer influences the strength and character of the higher-energy excitations. \nWhile in the 4T-F4TCNQ complex the first three bright excitations retain the same oscillator strength, in the presence of thiophene oligomers with six rings or more, higher-energy peaks gain significant spectral weight relative to the absorption onset and experience a sizable red-shift.\nAs a result, the optical density in the visible region is significantly enhanced upon increasing donor length.\nThe nature of the excitations changes accordingly. \nIn 4T-F4TCNQ the first three excitations have the same Frenkel-like character within the CTC.\nConversely, in the complexes formed by longer nT chains the second bright peak corresponds to a charge-transfer-like excitation and the third one to a Frenkel exciton with increasing localization on the donor as the size of the latter grows.\n\nOur work and the rationale elaborated therein offer important insight into the electronic structure and optical excitations of charge-transfer complexes.\nWhile doping in organic semiconductors remains a debated topic, our findings unravel fundamental structure-property relations in the presence of partial charge transfer, which is one of the leading doping mechanisms in molecular materials.\nAs such, our results represent an important step forward towards a deeper understanding of the opto-electronic properties of doped organic semiconductors.\n\n\\begin{acknowledgement}\nInspiring discussions with Arrigo Calzolari, Andreas Optiz, and Dieter Neher are gratefully acknowledged. \nThis work is funded by the Deutsche Forschungsgemeinschaft (DFG, German Research Foundation) - Projektnummer 182087777 - SFB 951 and HE 5866\/2-1.\n\\end{acknowledgement}\n\n\n\\begin{suppinfo}\nAdditional details about the structural, electronic, and optical properties of the CTCs are reported in the Supporting Information.\n\\end{suppinfo}\n\\newpage\n\n\\section*{Graphical TOC Entry}\n\\begin{figure}%\n\\centering\n\\includegraphics[height=4.45cm]{TOC_NEW_C.png}\n\\end{figure}\n\\newpage\n\\providecommand{\\latin}[1]{#1}\n\\providecommand*\\mcitethebibliography{\\thebibliography}\n\\csname @ifundefined\\endcsname{endmcitethebibliography}\n {\\let\\endmcitethebibliography\\endthebibliography}{}\n\\begin{mcitethebibliography}{45}\n\\providecommand*\\natexlab[1]{#1}\n\\providecommand*\\mciteSetBstSublistMode[1]{}\n\\providecommand*\\mciteSetBstMaxWidthForm[2]{}\n\\providecommand*\\mciteBstWouldAddEndPuncttrue\n {\\def\\unskip.}{\\unskip.}}\n\\providecommand*\\mciteBstWouldAddEndPunctfalse\n {\\let\\unskip.}\\relax}\n\\providecommand*\\mciteSetBstMidEndSepPunct[3]{}\n\\providecommand*\\mciteSetBstSublistLabelBeginEnd[3]{}\n\\providecommand*\\unskip.}{}\n\\mciteSetBstSublistMode{f}\n\\mciteSetBstMaxWidthForm{subitem}{(\\alph{mcitesubitemcount})}\n\\mciteSetBstSublistLabelBeginEnd\n {\\mcitemaxwidthsubitemform\\space}\n {\\relax}\n {\\relax}\n\n\\bibitem[Aziz \\latin{et~al.}(2007)Aziz, Vollmer, Eisebitt, Eberhardt, Pingel,\n Neher, and Koch]{aziz+07am}\nAziz,~E.; Vollmer,~A.; Eisebitt,~S.; Eberhardt,~W.; Pingel,~P.; Neher,~D.;\n Koch,~N. Localized Charge Transfer in a Molecularly Doped Conducting Polymer.\n \\emph{Adv.~Mater.~} \\textbf{2007}, \\emph{19}, 3257--3260\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Yim \\latin{et~al.}(2008)Yim, Whiting, Murphy, Halls, Burroughes,\n Friend, and Kim]{yim+08am}\nYim,~K.-H.; Whiting,~G.~L.; Murphy,~C.~E.; Halls,~J. J.~M.; Burroughes,~J.~H.;\n Friend,~R.~H.; Kim,~J.-S. Controlling Electrical Properties of Conjugated\n Polymers via a Solution-Based p-Type Doping. \\emph{Adv.~Mater.~}\n \\textbf{2008}, \\emph{20}, 3319--3324\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Pingel \\latin{et~al.}(2010)Pingel, Zhu, Park, Vogel, Janietz, Kim,\n Rabe, Br\\'{e}das, and Koch]{ping+10jpcl}\nPingel,~P.; Zhu,~L.; Park,~K.~S.; Vogel,~J.-O.; Janietz,~S.; Kim,~E.-G.;\n Rabe,~J.~P.; Br\\'{e}das,~J.-L.; Koch,~N. Charge-Transfer Localization in\n Molecularly Doped Thiophene-Based Donor Polymers. \\emph{J.~Phys.~Chem.~Lett.}\n \\textbf{2010}, \\emph{1}, 2037--2041\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Gao \\latin{et~al.}(2013)Gao, Roehling, Li, Guo, Moul\\'{e}, and\n Grey]{gao+13jmcc}\nGao,~J.; Roehling,~J.~D.; Li,~Y.; Guo,~H.; Moul\\'{e},~A.~J.; Grey,~J.~K. The\n Effect of 2,3,5,6-tetrafluoro-7,7,8,8-tetracyanoquinodimethane Charge\n Transfer Dopants on the Conformation and Aggregation of\n Poly(3-hexylthiophene). \\emph{J.~Mater.~Chem.~C} \\textbf{2013}, \\emph{1},\n 5638--5646\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Pingel and Neher(2013)Pingel, and Neher]{ping-nehe13prb}\nPingel,~P.; Neher,~D. Comprehensive Picture of $p$-Type Doping of P3HT with the\n Molecular Acceptor F${}_{4}$TCNQ. \\emph{Phys.~Rev.~B} \\textbf{2013},\n \\emph{87}, 115209\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Duong \\latin{et~al.}(2013)Duong, Wang, Antono, Toney, and\n Salleo]{duon+13oe}\nDuong,~D.~T.; Wang,~C.; Antono,~E.; Toney,~M.~F.; Salleo,~A. The Chemical and\n Structural Origin of Efficient p-Type Doping in P3HT. \\emph{Org.~Electron.~}\n \\textbf{2013}, \\emph{14}, 1330 -- 1336\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Salzmann \\latin{et~al.}(2012)Salzmann, Heimel, Duhm, Oehzelt, Pingel,\n George, Schnegg, Lips, Blum, Vollmer, and Koch]{salz+13prl}\nSalzmann,~I.; Heimel,~G.; Duhm,~S.; Oehzelt,~M.; Pingel,~P.; George,~B.~M.;\n Schnegg,~A.; Lips,~K.; Blum,~R.-P.; Vollmer,~A. \\latin{et~al.}\n Intermolecular Hybridization Governs Molecular Electrical Doping.\n \\emph{Phys.~Rev.~Lett.~} \\textbf{2012}, \\emph{108}, 035502\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[M\\'{e}ndez \\latin{et~al.}(2013)M\\'{e}ndez, Heimel, Opitz, Sauer,\n Barkowski, Oehzelt, Soeda, Okamoto, Takeya, Arlin, Balandier, Geerts, Koch,\n and Salzmann]{mend+13acie}\nM\\'{e}ndez,~H.; Heimel,~G.; Opitz,~A.; Sauer,~K.; Barkowski,~P.; Oehzelt,~M.;\n Soeda,~J.; Okamoto,~T.; Takeya,~J.; Arlin,~J.-B. \\latin{et~al.} Doping of\n Organic Semiconductors: Impact of Dopant Strength and Electronic Coupling.\n \\emph{Angew.~Chem.~Int.~Ed.} \\textbf{2013}, \\emph{52}, 7751--7755\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Yang \\latin{et~al.}(2014)Yang, Li, Duhm, Tang, Kera, and\n Ueno]{yang+14ami}\nYang,~J.; Li,~Y.; Duhm,~S.; Tang,~J.; Kera,~S.; Ueno,~N. Molecular\n Structure-Dependent Charge Injection and Doping Efficiencies of Organic\n Semiconductors: Impact of Side Chain Substitution. \\emph{Adv.~Mater.~}\n \\textbf{2014}, \\emph{1}, 1300128\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[M\\'{e}ndez \\latin{et~al.}(2015)M\\'{e}ndez, Heimel, Winkler, Frisch,\n Opitz, Sauer, Wegner, Oehzelt, R\\\"{o}thel, Duhm, T\\\"{o}bbens, Koch, and\n Salzmann]{mend+15ncom}\nM\\'{e}ndez,~H.; Heimel,~G.; Winkler,~S.; Frisch,~J.; Opitz,~A.; Sauer,~K.;\n Wegner,~B.; Oehzelt,~M.; R\\\"{o}thel,~C.; Duhm,~S. \\latin{et~al.}\n Charge-Transfer Crystallites as Molecular Electrical Dopants. \\emph{Nature\n Commun.} \\textbf{2015}, \\emph{6}, 8560\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Salzmann \\latin{et~al.}(2016)Salzmann, Heimel, Oehzelt, Winkler, and\n Koch]{salz+16acr}\nSalzmann,~I.; Heimel,~G.; Oehzelt,~M.; Winkler,~S.; Koch,~N. Molecular\n Electrical Doping of Organic Semiconductors: Fundamental Mechanisms and\n Emerging Dopant Design Rules. \\emph{Acc.~Chem.~Res.~} \\textbf{2016},\n \\emph{49}, 370--378\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Fuzell \\latin{et~al.}(2016)Fuzell, Jacobs, Ackling, Harrelson, Huang,\n Larsen, and Moul\\'e]{fuze+16jpcl}\nFuzell,~J.; Jacobs,~I.~E.; Ackling,~S.; Harrelson,~T.~F.; Huang,~D.~M.;\n Larsen,~D.; Moul\\'e,~A.~J. Optical Dedoping Mechanism for P3HT: F4TCNQ\n Mixtures. \\emph{J.~Phys.~Chem.~Lett.} \\textbf{2016}, \\emph{7},\n 4297--4303\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Jacobs and Moul\\'{e}(2017)Jacobs, and Moul\\'{e}]{jaco-moul17am}\nJacobs,~I.~E.; Moul\\'{e},~A.~J. Controlling Molecular Doping in Organic\n Semiconductors. \\emph{Adv.~Mater.~} \\textbf{2017}, \\emph{29}, 1703063\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Duva \\latin{et~al.}(2018)Duva, Pithan, Zeiser, Reisz, Dieterle,\n Hofferberth, Beyer, Bogula, Opitz, Kowarik, Hinderhofer, Gerlach, and\n Schreiber]{duva+18jpcc}\nDuva,~G.; Pithan,~L.; Zeiser,~C.; Reisz,~B.; Dieterle,~J.; Hofferberth,~B.;\n Beyer,~P.; Bogula,~L.; Opitz,~A.; Kowarik,~S. \\latin{et~al.} Thin-Film\n Texture and Optical Properties of Donor\/Acceptor Complexes.\n Diindenoperylene\/F6TCNNQ vs Alpha-Sexithiophene\/F6TCNNQ.\n \\emph{J.~Phys.~Chem.~C} \\textbf{2018}, \\emph{122}, 18705--18714\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Beyer \\latin{et~al.}(2019)Beyer, Pham, Peter, Koch, Meister,\n Br\\\"utting, Grubert, Hecht, Nabok, Cocchi, Draxl, and Opitz]{beye+19cm}\nBeyer,~P.; Pham,~D.; Peter,~C.; Koch,~N.; Meister,~E.; Br\\\"utting,~W.;\n Grubert,~L.; Hecht,~S.; Nabok,~D.; Cocchi,~C. \\latin{et~al.} State of Matter\n Dependent Charge Transfer Interactions between Planar Molecules for Doping\n Applications. \\emph{Chem.~Mater.~} \\textbf{2019},\n \\emph{10.1021\/acs.chemmater.8b01447}\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Jacobs \\latin{et~al.}(2018)Jacobs, Cendra, Harrelson, Bedolla~Valdez,\n Faller, Salleo, and Moul\\'{e}]{jaco+18mh}\nJacobs,~I.~E.; Cendra,~C.; Harrelson,~T.~F.; Bedolla~Valdez,~Z.~I.; Faller,~R.;\n Salleo,~A.; Moul\\'{e},~A.~J. Polymorphism Controls the Degree of Charge\n Transfer in a Molecularly Doped Semiconducting Polymer. \\emph{Mater.~Horiz.~}\n \\textbf{2018}, \\emph{5}, 655--660\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Neelamraju \\latin{et~al.}(2018)Neelamraju, Watts, Pemberton, and\n Ratcliff]{neel+18jpcl}\nNeelamraju,~B.; Watts,~K.~E.; Pemberton,~J.~E.; Ratcliff,~E.~L. Correlation of\n Coexistent Charge Transfer States in F4TCNQ-Doped P3HT with Microstructure.\n \\emph{J.~Phys.~Chem.~Lett.} \\textbf{2018}, \\emph{9}, 6871--6877\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Karpov \\latin{et~al.}(2017)Karpov, Erdmann, Stamm, Lappan, Guskova,\n Malanin, Raguzin, Beryozkina, Bakulev, G\\\"{u}nther, Gemming, Seifert,\n Hambsch, Mannsfeld, Voit, and Kiriy]{karp+17mm}\nKarpov,~Y.; Erdmann,~T.; Stamm,~M.; Lappan,~U.; Guskova,~O.; Malanin,~M.;\n Raguzin,~I.; Beryozkina,~T.; Bakulev,~V.; G\\\"{u}nther,~F. \\latin{et~al.}\n Molecular Doping of a High Mobility Diketopyrrolopyrrole--Dithienylthieno [3,\n 2-b] thiophene Donor--Acceptor Copolymer with F6TCNNQ. \\emph{Macromolecules}\n \\textbf{2017}, \\emph{50}, 914--926\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Scholes \\latin{et~al.}(2017)Scholes, Yee, Lindemuth, Kang, Onorato,\n Ghosh, Luscombe, Spano, Tolbert, and Schwartz]{scho+17afm}\nScholes,~D.~T.; Yee,~P.~Y.; Lindemuth,~J.~R.; Kang,~H.; Onorato,~J.; Ghosh,~R.;\n Luscombe,~C.~K.; Spano,~F.~C.; Tolbert,~S.~H.; Schwartz,~B.~J. The Effects of\n Crystallinity on Charge Transport and the Structure of Sequentially Processed\n F4TCNQ-Doped Conjugated Polymer Films. \\emph{Adv.~Funct.~Mater.~}\n \\textbf{2017}, \\emph{27}, 1702654\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Blase and Attaccalite(2011)Blase, and Attaccalite]{blas-atta11apl}\nBlase,~X.; Attaccalite,~C. Charge-Transfer Excitations in Molecular\n Donor--Acceptor Complexes within the Many-Body Bethe-Salpeter Approach.\n \\emph{Appl.~Phys.~Lett.~} \\textbf{2011}, \\emph{99}, 171909\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Baumeier \\latin{et~al.}(2012)Baumeier, Andrienko, and\n Rohlfing]{baum+12jctc}\nBaumeier,~B.; Andrienko,~D.; Rohlfing,~M. Frenkel and Charge-Transfer\n Excitations in Donor--Acceptor Complexes from Many-Body Green's Functions\n Theory. \\emph{J.~Chem.~Theory.~Comput.~} \\textbf{2012}, \\emph{8},\n 2790--2795\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Faber \\latin{et~al.}(2014)Faber, Boulanger, Attaccalite, Duchemin, and\n Blase]{fabe+14ptrsa}\nFaber,~C.; Boulanger,~P.; Attaccalite,~C.; Duchemin,~I.; Blase,~X. Excited\n States Properties of Organic Molecules: From Density Functional Theory to the\n GW and Bethe--Salpeter Green's Function Formalisms. \\emph{Phil. Trans. R.\n Soc. A} \\textbf{2014}, \\emph{372}, 20130271\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Niedzialek \\latin{et~al.}(2015)Niedzialek, Duchemin, de~Queiroz,\n Osella, Rao, Friend, Blase, K{\\\"u}mmel, and Beljonne]{nied+15afm}\nNiedzialek,~D.; Duchemin,~I.; de~Queiroz,~T.~B.; Osella,~S.; Rao,~A.;\n Friend,~R.; Blase,~X.; K{\\\"u}mmel,~S.; Beljonne,~D. First Principles\n Calculations of Charge Transfer Excitations in Polymer--Fullerene Complexes:\n Influence of Excess Energy. \\emph{Adv.~Funct.~Mater.~} \\textbf{2015},\n \\emph{25}, 1972--1984\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Zhu \\latin{et~al.}(2011)Zhu, Kim, Yi, and Br\\'{e}das]{zhu+11cm}\nZhu,~L.; Kim,~E.-G.; Yi,~Y.; Br\\'{e}das,~J.-L. Charge Transfer in Molecular\n Complexes with 2,3,5,6-Tetrafluoro-7,7,8,8-tetracyanoquinodimethane\n (F4-TCNQ): A Density Functional Theory Study. \\emph{Chem.~Mater.~}\n \\textbf{2011}, \\emph{23}, 5149--5159\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Blum \\latin{et~al.}(2009)Blum, Gehrke, Hanke, Havu, Havu, Ren, Reuter,\n and Scheffler]{blum+09cpc}\nBlum,~V.; Gehrke,~R.; Hanke,~F.; Havu,~P.; Havu,~V.; Ren,~X.; Reuter,~K.;\n Scheffler,~M. Ab initio Molecular Simulations with Numeric Atom-Centered\n Orbitals. \\emph{Comput.~Phys.~Commun.~} \\textbf{2009}, \\emph{180}, 2175 --\n 2196\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Havu \\latin{et~al.}(2009)Havu, Blum, Havu, and Scheffler]{havu+09jcp}\nHavu,~V.; Blum,~V.; Havu,~P.; Scheffler,~M. Efficient O(N) Integration for\n All-Electron Electronic Etructure Calculation Using Numeric Basis Functions.\n \\emph{J.~Comp.~Phys.~} \\textbf{2009}, \\emph{228}, 8367 -- 8379\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Perdew \\latin{et~al.}(1996)Perdew, Burke, and Ernzerhof]{perd+96prl}\nPerdew,~J.~P.; Burke,~K.; Ernzerhof,~M. Generalized Gradient Approximation Made\n Simple. \\emph{Phys.~Rev.~Lett.~} \\textbf{1996}, \\emph{77}, 3865--3868\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Tkatchenko and Scheffler(2009)Tkatchenko, and\n Scheffler]{tkat-sche09prl}\nTkatchenko,~A.; Scheffler,~M. Accurate Molecular Van Der Waals Interactions\n from Ground-State Electron Density and Free-Atom Reference Data.\n \\emph{Phys.~Rev.~Lett.~} \\textbf{2009}, \\emph{102}, 073005\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Bruneval \\latin{et~al.}(2016)Bruneval, Rangel, Hamed, Shao, Yang, and\n Neaton]{brun+16cpc}\nBruneval,~F.; Rangel,~T.; Hamed,~S.~M.; Shao,~M.; Yang,~C.; Neaton,~J.~B. molgw\n 1: Many-Body Perturbation Theory Software for Atoms, Molecules, and Clusters.\n \\emph{Comput.~Phys.~Commun.~} \\textbf{2016}, \\emph{208}, 149 -- 161\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Bruneval(2012)]{brun12jcp}\nBruneval,~F. Ionization Energy of Atoms Obtained from GW Self-Energy or from\n Random Phase Approximation Total Energies. \\emph{J.~Chem.~Phys.~}\n \\textbf{2012}, \\emph{136}, 194107\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Weigend \\latin{et~al.}(2002)Weigend, K\\\"{o}hn, and\n H\\\"{a}ttig]{weig+02}\nWeigend,~F.; K\\\"{o}hn,~A.; H\\\"{a}ttig,~C. Efficient Use of the Correlation\n Consistent Basis Sets in Resolution of the Identity MP2 Calculations.\n \\emph{J.~Chem.~Phys.~} \\textbf{2002}, \\emph{116}, 3175--3183\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Yanai \\latin{et~al.}(2004)Yanai, Tew, and Handy]{yana+04cpl}\nYanai,~T.; Tew,~D.~P.; Handy,~N.~C. A New Hybrid Exchange-Correlation\n Functional Using the Coulomb-Attenuating Method (CAM-B3LYP).\n \\emph{Chem.~Phys.~Lett.~} \\textbf{2004}, \\emph{393}, 51 -- 57\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Faber \\latin{et~al.}(2013)Faber, Boulanger, Duchemin, Attaccalite, and\n Blase]{fabe+13jcp}\nFaber,~C.; Boulanger,~P.; Duchemin,~I.; Attaccalite,~C.; Blase,~X. Many-body\n Green's Function GW and Bethe-Salpeter Study of the Optical Excitations in a\n Paradigmatic Model Dipeptide. \\emph{J.~Chem.~Phys.~} \\textbf{2013},\n \\emph{139}, 194308\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Bruneval and Marques(2013)Bruneval, and Marques]{brun+13jctc}\nBruneval,~F.; Marques,~M. A.~L. Benchmarking the Starting Points of the GW\n Approximation for Molecules. \\emph{J.~Chem.~Theory.~Comput.~} \\textbf{2013},\n \\emph{9}, 324--329\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Cocchi and Draxl(2015)Cocchi, and Draxl]{cocc-drax15prb}\nCocchi,~C.; Draxl,~C. Optical Spectra from Molecules to Crystals: Insight from\n Many-Body Perturbation Theory. \\emph{Phys.~Rev.~B} \\textbf{2015}, \\emph{92},\n 205126\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Palummo \\latin{et~al.}(2009)Palummo, Hogan, Sottile, Bagal{\\'a}, and\n Rubio]{palu+09jcp}\nPalummo,~M.; Hogan,~C.; Sottile,~F.; Bagal{\\'a},~P.; Rubio,~A. Ab initio\n Electronic and Optical Spectra of Free-Base Porphyrins: The Role of\n Electronic Correlation. \\emph{J.~Chem.~Phys.~} \\textbf{2009}, \\emph{131},\n 08B607\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Baumeier \\latin{et~al.}(2012)Baumeier, Andrienko, Ma, and\n Rohlfing]{baum+12jctc1}\nBaumeier,~B.; Andrienko,~D.; Ma,~Y.; Rohlfing,~M. Excited states of\n Dicyanovinyl-Substituted Oligothiophenes from Many-Body Green's Functions\n Theory. \\emph{J.~Chem.~Theory.~Comput.~} \\textbf{2012}, \\emph{8},\n 997--1002\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Rangel \\latin{et~al.}(2017)Rangel, Hamed, Bruneval, and\n Neaton]{rang+17jcp}\nRangel,~T.; Hamed,~S.~M.; Bruneval,~F.; Neaton,~J.~B. An Assessment of\n Low-Lying Excitation Energies and Triplet Instabilities of Organic Molecules\n with an Ab initio Bethe-Salpeter Equation Approach and the Tamm-Dancoff\n Approximation. \\emph{J.~Chem.~Phys.~} \\textbf{2017}, \\emph{146}, 194108\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Cocchi \\latin{et~al.}(2011)Cocchi, Prezzi, Ruini, Caldas, and\n Molinari]{cocc+11jpcl}\nCocchi,~C.; Prezzi,~D.; Ruini,~A.; Caldas,~M.~J.; Molinari,~E. Optical\n Properties and Charge-Transfer Excitations in Edge-Functionalized\n All-Graphene Nanojunctions. \\emph{J.~Phys.~Chem.~Lett.} \\textbf{2011},\n \\emph{2}, 1315--1319\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[De~Corato \\latin{et~al.}(2014)De~Corato, Cocchi, Prezzi, Caldas,\n Molinari, and Ruini]{deco+14jpcc}\nDe~Corato,~M.; Cocchi,~C.; Prezzi,~D.; Caldas,~M.~J.; Molinari,~E.; Ruini,~A.\n Optical Properties of Bilayer Graphene Nanoflakes. \\emph{J.~Phys.~Chem.~C}\n \\textbf{2014}, \\emph{118}, 23219--23225\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Fabiano \\latin{et~al.}(2005)Fabiano, Della~Sala, Cingolani, Weimer,\n and G\\\"{o}rling]{fabi+05jpca}\nFabiano,~E.; Della~Sala,~F.; Cingolani,~R.; Weimer,~M.; G\\\"{o}rling,~A.\n Theoretical Study of Singlet and Triplet Excitation Energies in\n Oligothiophenes. \\emph{J.~Phys.~Chem.~A} \\textbf{2005}, \\emph{109},\n 3078--3085\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Sun and Autschbach(2014)Sun, and Autschbach]{sun-auts14jctc}\nSun,~H.; Autschbach,~J. Electronic Energy Gaps for $\\pi$-conjugated Oligomers\n and Polymers Calculated with Density Functional Theory.\n \\emph{J.~Chem.~Theory.~Comput.~} \\textbf{2014}, \\emph{10}, 1035--1047\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Hogan \\latin{et~al.}(2013)Hogan, Palummo, Gierschner, and\n Rubio]{hoga+13jcp}\nHogan,~C.; Palummo,~M.; Gierschner,~J.; Rubio,~A. Correlation Effects in the\n Optical Spectra of Porphyrin Oligomer Chains: Exciton Confinement and Length\n Dependence. \\emph{J.~Chem.~Phys.~} \\textbf{2013}, \\emph{138}, 01B607\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\bibitem[Li \\latin{et~al.}(2017)Li, D'Avino, Pershin, Jacquemin, Duchemin,\n Beljonne, and Blase]{li+17prm}\nLi,~J.; D'Avino,~G.; Pershin,~A.; Jacquemin,~D.; Duchemin,~I.; Beljonne,~D.;\n Blase,~X. Correlated Electron-Hole Mechanism for Molecular Doping in Organic\n Semiconductors. \\textbf{2017}, \\emph{1}, 025602\\relax\n\\mciteBstWouldAddEndPuncttrue\n\\mciteSetBstMidEndSepPunct{\\mcitedefaultmidpunct}\n{\\mcitedefaultendpunct}{\\mcitedefaultseppunct}\\relax\n\\unskip.}\n\\end{mcitethebibliography}\n\n\n\n\\end{document}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nDespite the fact that ultra-high energy (UHE) cosmic rays are known\nfor decades, their sources and the acceleration mechanisms are\nstill under debate. Sources are only detectable by $\\gamma$-rays produced\nin interactions near the sources.\nIn the very-high energy (VHE) range near 1\\,TeV \nmore and more $\\gamma$-ray sources are revealed by the imaging Cherenkov technique.\nOn the contrary, no clear source detections have been made in the UHE domain\nabove about 100\\,TeV, except perhaps a few episodic cases. \n\nMainly for reasons of the required power, the dominant sources of\ncosmic rays up to about 100\\,TeV and probably up to the {\\em knee}\nof the cosmic-ray energy spectrum at a few $10^{15}$\\,eV are believed\nto be supernova remnants in the Sedov phase. The change of the spectrum\nnear the knee presumably reflects a change in the origin and the\ntakeover of another, yet unclear type of sources at energies above\nthe knee. A change in the cosmic-ray propagation with a decreasing\nGalactic containment has also been considered. In either case,\nthe change in the slope of the spectrum should be accompanied by a change \nin the mass composition of cosmic rays. In the case of a change of\nsources across the knee, the composition could change dramatically.\nA much less spectacular change of the composition is expected\nin the case of decreasing containment.\n\nDirect measurements well below the knee\n\\cite{Asakimori-1993ab,Ichimura-1993b}\nshow indeed a substantial change in the mass composition\nalready at energies around 100\\,TeV. In particular, the fraction\nof protons seems to diminish with increasing energy.\nDue to their small collection\nareas the balloon-borne direct experiments run out of statistics\nabove several hundred TeV. Indirect methods using ground-based experiments\nare, so far, not able to classify individual cosmic rays unambiguously\nby their mass. Such methods are more appropriate for evaluating some\naverage mass number. Most notably, the results of the Fly's Eye\ngroup \\cite{Gaisser-1993} indicate a rather heavy composition well above \nthe knee. If taken at face value, their results represent a composition of\nmainly very heavy nuclei, like iron, at $10^{17}$\\,eV.\n\nExperiments measuring the composition right at the knee of the\nspectrum obtained either ambiguous or even conflicting results.\nResults with no significant change\n\\cite{Zhu-1990,Ahlen-1992b} or a slight increase of heavy elements\n\\cite{Khristiansen-1994} have been reported. Other \ngroups found more significant increases\nof heavy elements \\cite{Ren-1988a,Freudenreich-1990,Mitsui-1995} \nor, on the other hand, predominantly protons \\cite{Cebula-1990}.\nFurther measurements are needed, if the cosmic-ray composition\nshould help to resolve the question where cosmic rays are accelerated.\n\nThe method proposed in this paper should be relatively easy to implement\nat sites where an air-shower array with an angular resolution below\none degree already exists. The method is mainly a measurement of the\nlongitudinal shower development by the angles of muons with respect to\nthe shower axis. It does not require to measure many muons in a single \nshower because it uses only the {\\em inclusive}\nangular distributions, nor is accurate timing required.\nIndeed, if several muons are measured in one event, they are treated\nseparately.\n\nNot only the average muon angles but also the average electron angles \nwith respect to the shower axis are sensitive to the primary\ncomposition. Because the detector response \nis, in general, better understood for muon tracks than for electron tracks,\nthis paper focuses mainly on the muons.\nUsing the tracking detectors and the air-shower array, the\nmethod can be supplemented by traditional methods like the\naverage $\\mu\/e$ ratio or the muon lateral distribution.\nAlthough one can think of complex experiments dedicated\nto measuring the cosmic-ray composition, like KASKADE \\cite{Rebel-1993},\nwhere more pieces of information are collected, the\npurpose of this paper is to demonstrate that even \nthe muon angles alone provide significant information about\nthe composition.\n\nThe proposed method relies on shower simulations to obtain\nan absolute value of an average mass number, just like\nother indirect methods. If suitable tracking detectors are used,\nit has the advantage that essentially\nall required detector parameters for the simulation can be\nobtained from the measured data as a function of shower size.\nTherefore, it can be particularly sensitive to changes in the composition\nacross the measured shower size spectrum, even if a comparison\nwith simulations using different interaction models may yield\ndifferent absolute values.\n\nAn analysis of data taken with ten CRT detectors \n-- tracking detectors of 2.5\\,m$^2$ sensitive area each \\cite{CRT-NIM-1} --\nand the HEGRA air-shower array \\cite{Fonseca-1995} on La Palma is\nin progress \\cite{future-CRT-results}. In the data of CRT and HEGRA\nboth muon and electron tracks are analysed. Although the intention of\nthis paper is to promote the method to other air-shower experiments,\nthe simulations shown in this paper are specific to the combination of the\nCRT and HEGRA experiments. These simulations take into account the\nresponse of both components -- the tracking detectors and the\nair-shower array -- in much detail. After an outline of the method,\nthe specific simulations are described to demonstrate that the\nmethod can be very well applied with existing detector technology.\nExperimental requirements for application of the method and\nlimitations by shower simulations are shown in a more general context,\nnot specific to CRT and HEGRA.\n\n\n\\section{Outline of the method}\n\nFirst measurements of the angles of muons with respect to the\ncorresponding air showers or, in alternative terms, of the apparent\nheight of origin of the muons were done in the 1960s \\cite{Earnshaw-1973}.\nAt that time neither the angular resolutions of air-shower arrays\nnor the accuracy of shower simulations were appropriate to use\nsuch measurements for investigating the mass composition of\nprimary cosmic rays. For an early comparison of measured\nand simulated height of muon origin see \\cite{Gaisser-1978}.\n\nThe proposed method is based on the fact that showers initiated\nby heavy primaries, like iron nuclei, develop on average\nat smaller atmospheric depths, i.e.\\ higher altitudes, than proton showers.\nMuons from showers of heavy primaries arrive, therefore, at smaller\nangles with respect to the shower axis for any given core distance.\nThe proposed method is not intended to measure individual primaries\nbut to measure accurately an average number of the composition, like \n$\\langle\\ln A\\rangle$, the average logarithm of the primary mass numbers.\nMeasuring many tracks in a single shower is, therefore, not necessary\nand only helps to improve statistics.\n\n\\begin{figure}[htbp]\n\\epsfxsize=0.6\\textwidth\n\\hc{\\epsffile{geo_rad_tang.eps}}\n\\caption[Geometry of radial and tangential angles]{The\ngeometry of radial and tangential angles between the\nshower axis, e.g.\\ as reconstructed from air-shower data,\nand the measured track of a particle.}\n\\label{fig:geometry}\n\\end{figure}\n\nThe angles of muons or other particles with respect to the shower axis are \nbest expressed in terms of their {\\em radial angle\\\/} and {\\em tangential angle\\\/}\ncomponents. The radial angle is the projection onto the plane defined\nby the geometry of the shower axis and the location of the muon detector.\nThroughout this paper, a negative radial angle is assigned to\na particle flying away from the shower axis. In fact, at large\ncore distances most particles have such a negative radial angle.\nThe tangential angle (also called the transverse angle) is the projection\nonto the plane perpendicular to the `radial' plane and parallel to the\nshower axis (see figure~\\ref{fig:geometry}). \nFor almost vertical\nshowers a more convenient definition with projections into vertical\nplanes defined only by the positions of the shower core and the\ndetector is essentially equivalent. \n\nThe radial projection can be\nused to derive an {\\em apparent height of origin\\\/} of a particle\n(figure~\\ref{fig:geo-height}). Although radial angle and apparent\nheight of origin are almost equivalent, the radial angle is the\nmore robust variable.\nThis is due to the fact that\na small scattering or experimental smearing around zero radial\nangle corresponds to $\\pm\\infty$ apparent height.\nIn addition, the derived height for each muon depends directly on the\nresolution of the core position. In particular, the last point turns out \nto be a major drawback of an analysis in terms of the apparent height.\nThe core position resolution is usually not very precisely known\nas a function of shower size, zenith angle, and position within the\narray. An analysis in terms of the apparent height is, for this\nreason, limited to core distances much larger than the core position\nresolution.\n\n\\begin{figure}[htbp]\n\\epsfxsize=0.6\\textwidth\n\\hc{\\epsffile{geo_height.eps}}\n\\caption[Apparent height of origin]{An apparent\nheight of origin can be derived from the track projected onto\nthe radial plane.}\n\\label{fig:geo-height}\n\\end{figure}\n\nThe distribution of tangential angles (see figure \\ref{fig:tang+rad})\nis symmetric and its width is\ndominated by multiple scattering and the lateral distribution of the\nparent particle generation. For the muons in UHE air showers the\nintrinsic distribution is very narrow, apart from tails due\nto some low-energy muons. This distribution shows\nlittle difference between proton- and iron-initiated showers,\neven for an ideal angular resolution.\nUnder certain experimental constraints the tangential angle\ncan be used to measure the combined angular resolution of\nair-shower array and tracking detectors.\n\n\\begin{figure}[htbp]\n\\epsfxsize=0.8\\textwidth\n\\hc{\\epsffile{tan-rad-ideal.eps}}\n\\caption[Distributions of tangential and radial muon angles]{Distributions\nof tangential and radial angles of shower muons in simulated \nproton and iron showers, respectively,\nwith $N_e>10000$, at core distances of 200--250~m. \nPerfect angular resolution but a realistic detection\nefficiency as appropriate for CRT detectors is assumed.}\n\\label{fig:tang+rad}\n\\end{figure}\n\nThe radial angle, representing the longitudinal shower development,\nis the quantity sensitive to the composition. The {\\em average}\n(which may be mean, median, or mode) radial angle depends in a\ncharacteristic way on the core distance. \nTo compare the simulated average radial-angle curves\nwith measurements requires to know the combined angular resolution of the\nair-shower array and tracking detectors and to some extent also the\ncore position resolution of the array. However, the average radial angle\nat any given core distance changes very little with shower size or zenith angle.\nThe figures in this paper are intended as examples. They are based on\nsimulations using the CORSIKA code (see section \\ref{sect:shower+det-sim}).\nThe simulations assume an experimental resolution as observed\nfor CRT detectors at the HEGRA array \\cite{CRT-NIM-2}. \nSee figure~\\ref{fig:rad-ang-primaries}\nfor the median radial angles of different primary types.\nAt any core distance the average radial angle is almost a linear function\nof $\\langle\\ln A\\rangle$, where $A$ is the mass number of the primary.\n\n\\begin{figure}[htbp]\n\\epsfxsize=0.8\\textwidth\n\\hc{\\epsffile{rad-ang-primaries.eps}}\n\\caption[Radial angle of muons for different primaries]{Median radial angle of \nmuons in showers initiated by different primaries, for a minimum shower size of \n10$^4$ and angular resolutions and detector responses as \nappropriate for CRT detectors and the HEGRA scintillator array \n(see section \\ref{sect:shower+det-sim}).}\n\\label{fig:rad-ang-primaries}\n\\end{figure}\n\nAlthough a similar difference between showers initiated by protons or\niron nuclei is also present in the radial-angle curves of electrons,\nthe electron component is less usefull for composition studies than muons.\nFirst of all, a very detailed simulation of the behaviour of\nthe tracking detectors to electrons, gammas, muons, and hadrons \nis required to compare with experimental data. \nFor muons the detector is generally much better understood,\nalthough punch-throughs of other particles have to be considered.\n\nThe second reason in favour of muons is based on the shower selection\nintroduced by the air-shower array. To achieve the same shower size\n($N_e$) at ground level, showers initiated by heavy primaries require\na larger primary energy than proton showers. Due to the steep spectrum,\nconventional air-shower arrays will, therefore, mainly see proton showers\nat any shower size, even if protons account for only half of the\nmass composition at a given energy per particle. Muons can almost\ncompensate for this bias because, at the same shower size, a proton\nshower contains less muons than a shower of a heavy primary.\n\nBecause of the compensation of the selection bias by the larger number\nof muons for heavy primaries the quantity\n\\begin{equation}\n\\Lambda_\\mu = \\frac{1}{\\sum w_i}\\sum_i \\,w_i\\,\\frac{\\langle\\alpha_i\\rangle -\n \\langle\\alpha_{i,{\\rm p}}\\rangle}{\\langle\\alpha_{i,{\\rm Fe}}\\rangle -\n \\langle\\alpha_{i,{\\rm p}}\\rangle}\n\\end{equation}\nis even for a mixed composition essentially proportional to $\\langle\\ln A\\rangle$:\n\\begin{equation}\n \\langle\\ln A \\rangle \\approx \\Lambda_\\mu\\,\\ln 56.\n\\label{eq:lnA-Lambda}\n\\end{equation}\nIn this context $\\langle\\alpha_i\\rangle$ is the measured (or simulated) \nmedian radial angle in the core distance interval $i$, \n$\\langle\\alpha_{i,{\\rm p}}\\rangle$ and $\\langle\\alpha_{i,{\\rm Fe}}\\rangle$\nare the expected median values for pure protons and pure iron \nnuclei from the simulations, and the weights $w_i$ take\nthe statistical accuracy of measured as well as simulated average values\ninto account. In the simplest possible formulation, ignoring statistical\nerrors of the simulation and correlations of errors in different\nradial bins and taking into account only $\\sigma(\\langle\\alpha_i\\rangle)$\nas the measured accuracy of $\\langle\\alpha_i\\rangle$,\n\\begin{equation}\n w_i = \\Bigl( \\sigma(\\langle\\alpha_i\\rangle) \\;\/\\; \n ({\\langle\\alpha_{i,{\\rm Fe}}\\rangle -\n \\langle\\alpha_{i,{\\rm p}}\\rangle}) \\Bigr)^{-2}.\n\\end{equation}\n\nFor compositions as reported by direct measurements, equation\n\\ref{eq:lnA-Lambda}\nholds to\nbetter than 0.05 at all investigated shower size intervals --\nat least for CRT and HEGRA to which the simulations\ncorrespond. Although this relation is quite coincidental,\nit is expected to hold as well for most other existing\nair-shower arrays in combination with tracking detectors\nof a low energy threshold for muons.\nIn the case of extreme compositions like\npure helium, however, the relation may be wrong by up to 0.26\nin terms of $\\langle\\ln A\\rangle$.\n\n\\begin{figure}[htbp]\n\\epsfxsize=0.8\\textwidth\n\\hc{\\epsffile{rad-ang-composite.eps}}\n\\caption[Radial angle of muons for different compositions]{Median radial \nangle of muons in showers\nfor different energy-independent model compositions (pure protons, pure iron,\nand two compositions as measured by the JACEE collaboration \n\\cite{Asakimori-1993ab}).\nSimulations for $N_e>10^4$, as in figure \\ref{fig:rad-ang-primaries}.}\n\\label{fig:mu-rad-ang-composite}\n\\end{figure}\n\n\\begin{figure}[htbp]\n\\epsfxsize=0.8\\textwidth\n\\hc{\\epsffile{elec-rad-ang-composite.eps}}\n\\caption[Radial angle of electrons for different compositions]{Simulated \nmedian radial angle of particles identified as electrons,\nfor the same model compositions as in figure~\\ref{fig:mu-rad-ang-composite}.\n$N_e>10^4$, as in figure \\ref{fig:rad-ang-primaries}.}\n\\label{fig:elec-rad-ang-composite}\n\\end{figure}\n\nFigure \\ref{fig:mu-rad-ang-composite} shows that the median (or mean) radial\nangle of muons is indeed sensitive to changes in the composition\nas indicated by direct measurements \\cite{Asakimori-1993ab} below\nthe knee. Even the electron radial angles can be used for that\npurpose (see figure \\ref{fig:elec-rad-ang-composite}) but\na more accurate knowledge of the detector response is required in that case.\nIn addition, $\\Lambda_e$ (defined in analogy to $\\Lambda_\\mu$), is not\nproportional to $\\langle\\ln A\\rangle$ but is more sensitive to the\nfraction of light nuclei than to heavy nuclei.\n\n\n\\section{Shower and detector simulations}\n\\label{sect:shower+det-sim}\n\nThe scope of this paper is beyond the measurements\ncarried out with CRT detectors and the HEGRA array.\nAlthough the simulations presented here serve only as examples,\nexperimental effects of CRT detectors and the\nHEG\\-RA scintillator array are carefully taken into account. It may\nbe instructive to see which of these effects turn out to be\nthe most relevant ones.\n\nThe CORSIKA simulation code \\cite{CORSIKA-1992,Knapp-1993}\nwas used in version 4.068 for all shower simulations presented in this paper.\nExcept for tests of the accuracy of the shower simulations (see section\n\\ref{sect:shower-accuracy}) the VENUS hadron-nucleus and nucleus-nucleus\ninteraction code \\cite{Werner-1993} was used for high-energy\nhadronic interactions, GHEISHA \\cite{Fesefeldt-1985} for \nhadronic interactions below 80~GeV, and EGS4 \\cite{Nelson-1985} for\nthe electromagnetic component -- all in the framework of CORSIKA.\nThe simulations take also the geomagnetic field at La Palma into account.\nA total of 5580 showers (1860 proton showers and 930 showers\nfor each of the other four types of primaries, He, N, Mg, and Fe) were generated \nisotropically in the zenith angle range 0$^\\circ$ to 32$^\\circ$.\nThe energy ranges were selected to cover approximately the same\nshower sizes, from 20~TeV to 2~PeV for protons increasing to\n40~TeV to 4~PeV for iron nuclei. Showers were generated with a\n$E^{-1.7}$ differential energy spectrum to have sufficient\nnumbers of showers at all energies and weighted by $E^{-1}$ for the\nanalysis. The analysis was restricted to shower sizes where\ncontributions from showers outside the simulated energy ranges\nare negligible. \n\nConcerning the CRT detectors \\cite{CRT-NIM-1}, it should be \nmentioned that muons are identified as pairs of tracks in two\ndrift chambers, one above and one below a 10~cm thick iron plate.\nThe measured angles between the two tracks of an identified muon are\nrequired to be less than 2.5$^\\circ$ in two perpendicular projections. \nElectrons are identified as non-muon tracks in the upper drift chamber.\nResults of very detailed simulations of the response of CRT detectors to various\ntypes of particles (electrons, gammas, muons, pions, and protons) -- as\noutlined in \\cite{CRT-NIM-2} and described in detail in \\cite{Zink-1995} --\nand the measured detector performance \\cite{CRT-NIM-2}\nwere carefully parametrized. This includes detection efficiencies\nand angular resolutions as functions of particle energies, zenith angles,\nand track densities -- independent\nfor each particle type and for the identification as an electron track \nand as a muon track. As a consequence, the simulations include\nnot only punch-through of electrons but also of hadrons.\nIt should be noted that the energy dependence for\ndetecting genuine muons is fully described by multiple\nscattering and energy loss in the iron plate \\cite{CRT-NIM-2}.\nCRT detectors are mainly sensitive to muons above about\n1~GeV. Electrons, on the other hand, are detected above some 10~MeV.\nThe effect of multiple scattering in the detector container\nand the iron plate are included in the angular resolutions.\n\nFor a detailed simulation of the HEGRA array see \\cite{Martinez-1995}.\nFor the simulations presented here, a much simpler approach is\nsufficient because the simulations are restricted to shower\nsizes where the HEGRA array is fully efficient, independent of\nshower age, core position, and zenith angle.\nThe resolutions for shower sizes and core positions as obtained\nby the detailed simulations were parametrized. \nThe HEGRA angular resolution is implemented in the simulations\nin such a way that the combined\nangular resolution of CRT and HEGRA -- as measured by the muon\ntangential angles as a function of shower size and zenith angle \n\\cite{CRT-NIM-2} -- are fully reproduced.\n\nIt may be instructive to note that the median radial-angle curves\npresented in this paper are not significantly altered by\nmodifying the assumed angular resolution\nwithin the experimental limits. Modifying the assumed\ncore position resolution has an impact on the median radial\nangles only within a few ten meters from the core. \nPunch-throughs, although included in some detail, turned out\nto be of minor importance with CRT detectors, due to their\nexcellent punch-through rejection. The most important\nexperimental parameters in the simulations are the energy \ndependence of the detection efficiencies and the combined angular resolution\nof tracking detectors and air-shower array. In the case of CRT detectors\nand the HEGRA array both are very well known for muons and even\nsufficiently well known for electrons. \n\nIt appears very reasonable that the response of other types of \ntracking detectors and other air-shower arrays can also be \nunderstood well enough to take advantage of the radial-angle method. \nHowever, a few experimental requirements as outlined in section\n\\ref{sect:requirements} should be observed for that purpose.\n\n\n\\section{Accuracy of shower simulations}\n\\label{sect:shower-accuracy}\n\nApart from the detector response, a major uncertainty in the\ninterpretation of any indirect measurement of the\ncosmic-ray mass composition is due to the limited accuracy of\nthe shower simulations. Despite much progress in recent\nyears, the nucleus-nucleus and hadron-nucleus interaction models in \nair-shower simulations lead to non-negligible theoretical\nuncertainties. The CORSIKA simulation code\nattempts a particularly detailed and accurate shower simulation. \nIn the version 4.068, as used for this paper, it offers several\nindependent interaction models.\n\n\\begin{figure}[htbp]\n\\epsfxsize=0.8\\textwidth\n\\hc{\\epsffile{rad-ang-syserr.eps}}\n\\caption[Radial angle of muons for interaction models]{Median radial \nangle of muons in proton and iron showers with different interaction\nmodels (10000\\,$<$\\,$N_e$\\,$<$\\,20000, with fixed height of first interaction\nto keep shower-to-shower fluctuations small). 1: protons with\nVENUS\/GHEISHA, 2: protons with DPM, 3: iron with VENUS\/GHEISHA,\n4(5): iron with DPM and no (full) fragmentation of the spectators\nin projectile nuclei.}\n\\label{fig:rad-ang-syserr}\n\\end{figure}\n\nResults for different\ninteraction models, using the VENUS and GHEI\\-SHA \noptions on one side and DPM on the other side, and for\ndifferent fragmentation (evaporation) of the projectile nuclei have\nbeen compared in the energy range $10^{14}$--$10^{15}$\\,eV.\nSystematic differences in the average shower sizes of up\nto 20\\% are found and also differences in the\nshower fluctuations. If the median radial-angle curves are compared\nin the same ranges of shower sizes, the differences between models are\nvery small within core distances up to about 200\\,m, increasing to\nabout 10\\% of the proton-iron difference at a core distance of 400\\,m \n(figure~\\ref{fig:rad-ang-syserr}).\nThat increase of systematic uncertainties is due to the fact that \nhard interactions resulting in high-$p_t$ secondaries are not yet\nincluded very accurately.\nWith all interaction models the shower-size dependence of the\nradial-angle curve is small. A 30\\% systematic error in the\nrelation of simulated to measured shower sizes would result\nonly in a systematic error of the median radial angles of 3\\% of\nthe proton-iron difference.\nMoreover, this dependence is essentially \nthe same for different models.\n\nAnother uncertainty in the longitudinal shower development is\ndue to the sparse particle physics data in the very forward rapidity \nregion at high energies. Until more such data is available and \nincorporated into the interaction models it seems desirable to have also\nindirect composition measurements in an energy range overlapping with \ndirect measurements. When searching only for a change of the\ncomposition near the knee of the cosmic-ray spectrum, the poor\nknowledge of the very forward region should not be a limitation.\nDifferent parametrizations of parton densities which pose problems\nfor shower simulations at extremely high energies \\cite{Capdevielle-1995}\nare not of concern near $10^{15}$\\,eV primary energy.\n\n\n\n\\section{Experimental requirements}\n\\label{sect:requirements}\n\nTo exploit the described method for composition studies several\nexperimental requirements should be met. These are requirements for\nthe {\\em particle identification} with a tracking detector and \nthe {\\em resolution} of angles and positions.\n\nBecause the particles of the electromagnetic shower component, \n$e^\\pm$ and $\\gamma$, have angular distributions much different from\nthose of muons, a good electron rejection is important. Hadrons\nonly matter very close to the shower core. Due to the broad\nmuon lateral distribution it is generally not such a problem\nto identify muons at large core distances but more of a problem\nto reject random, non-shower muons. This can be easily solved\nby a timing information of a few hundred nanoseconds accuracy\nand by the angular correlation of shower muons with the shower axis.\nAnother requirement for the\nmuon identification is due to the fact that the radial angles\nare well correlated with the muon momentum, with higher-energy\nmuons usually coming from further up in the atmosphere and being\nbetter aligned with the shower axis. The muon identification should,\ntherefore, be well understood also as a function of momentum and\nzenith angle.\nFor a high muon-energy threshold, a very good angular resolution\nis required, while for a low threshold the multiple scattering\nand magnetic deflection of muons cannot be neglected.\n\nThe combined angular resolution of air-shower array\nand tracking detector should be better than about 1$^\\circ$\nin each projection but has to be known as a function of shower size\nto compare measurements with simulations. Definitely sufficient would be\na combined resolution of some 0.5$^\\circ$ if the muon-energy \nthreshold is as low as about 1~GeV.\nThis would be difficult to achieve with a tracking detector\nentirely below thick shielding material which also scatters the muons.\nA good uniformity of the resolution over azimuth angles would\nbe of advantage for the important detector calibration. Only with a\nuniform resolution can the angular resolution be derived from\nthe tangential angle distribution of the same data set which is\nused to compare the radial angles with simulations. Otherwise,\nshower-size dependent systematic errors can arise. \n\nThe alignment of the tracking detectors with respect to the\nreference frame as defined by showers reconstructed from array data\nshould be known to better than about 0.1$^\\circ$. This may require some\ncareful calibrations but is certainly feasible. \n\nThe core position resolution which also enters into the simulations \nhas to be determined from the air-shower array alone. Keeping \nselection effects of the air-shower array as small as possible\ndemands to select only showers well contained in the array\nand with shower sizes well above the trigger and reconstruction\nthresholds. For many of the existing air-shower arrays that\nwould correspond to threshold energies of some 80--200\\,TeV,\njust below where direct measurements are running out of statistics.\nReasonable errors in the shower size and energy calibrations of the array \ndo not compromise the method as long as the same experimental resolution\nis used in the simulations as seen in the data to be compared with\nthese simulations.\n\nWith the 25~m$^2$ sensitive area of the CRT detectors the\nstatistics with data of about one year would be the limit for \nprimary energies above several $10^{16}$~eV. A few hundred\nsquare meters area would be sufficient\nin order to achieve\nsome overlap with the Fly's Eye composition measurements.\n\nAll these requirements can be fulfilled with present\ntechnology for large-area tracking detectors and existing air-shower\narrays. Although simultaneous precise (nanosecond) timing\ninformation of the muons could supplement the described method\n(see for example \\cite{Danilova-1994}),\nprecise timing is not required and might be difficult to accomplish\nwith large-area detectors.\n\n\n\\section{Concluding remarks}\n\nThe average radial angle of particles are sensitive to the\ncosmic-ray mass composition. Muons have, among other particles in air\nshowers, many advantages. Their intrinsic angular distribution is\nvery narrow and, thus, the tangential-angle distribution can be\nused to calibrate the radial-angle resolution as a function of\nshower size. Muons can be distinguished from other particle types\nvery well. The larger number of muons in showers initiated by \nheavy primaries compensates for the $N_e$ shower selection bias\nwhich favours light primaries.\n\nLike any other\nindirect measurement of some average of the cosmic-ray composition,\nthat method requires comparison of measured data with simulations.\nCompared to many other indirect methods, the muon radial-angle method\nhas two major advantages: First, the most important detector effects\nto be included in the simulation can be measured very well and, second,\nsystematic errors, for example due to the interaction model,\ncan be easily checked by comparing measurements at low energies\n(a few hundred TeV) with simulations for directly measured\ncompositions as in \\cite{Asakimori-1993ab}. \n\nDespite some systematic uncertainties in the shower simulations\nan average mass number ($\\langle\\ln A\\rangle$) can be derived\nand compared with direct measurements well below the knee. \nRegardless of that, the muon radial-angle method can be very\nsensitive to changes in the composition. The radial-angle\nmethod would not require a very large experimental effort if\nexisting air-shower arrays are supplemented with suitable\ntracking detectors.\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\n\\section*{Ideas that we want to include}\n\n\\input{abstract}\\\\\n\\noindent\\textbf{Key Words}: BUGS\/JAGS; Capture-recapture; Data augmentation; Heterogeneity; Multinomial; Reversible jump MCMC.\n\n\\section{Introduction}\nInference about $N$, the unknown size of a population, based on capture-recapture models is of enormous interest, with ecology and epidemiology just two areas in which there are many applications. This is a specific case of a more general class of problems in which the index of a multinomial model is unknown. A sub-class of models of particular interest are the heterogeneity models in which capture probabilities vary among individuals. \\cite{Link2003} has shown that inference about $N$ within this class of models is sensitive to the choice of model for describing individual variation.\n\nA drawback of much available software for fitting $M_h$ is that analyses are restricted to default choices of model for the individual capture probabilities. Restricting analyses based on available computer code seems undesirable and an advantage of Bayesian approaches is that hierarchical modelling and data augmentation implemented via Markov chain Monte Carlo (MCMC) provide a natural framework for general model fitting.\n\nA commonly used dataset for illustrating model $M_h$ is the snowshoe hare data provided by \\cite{Otis1978} in which 68 individuals were captured in six days of trapping with encounter frequencies $f\\bm = (25, 22, 13, 5, 1, 2)'$ where $f_j \\; (i=1, \\ldots,k)$ is the number of hares caught $i$ times during the $k=6$ study. Define $N$ as the population size, $y_i$ as the number of times that individual $i$ was caught and $p_i$ as the corresponding capture probability. To model these using a program such as BUGS \\cite[]{Lunn2000} or JAGS \\cite[]{Plummer2003} we would like to be able to specify the model simply as\n\\[\n y_i \\sim Bin(k, p_i)\n\\]\nfor $i$ in $1:N$ (herein, unless we specify the two programs separately, we will use BUGS to refer to BUGS and JAGS jointly). To complete the model specification we would then specify a suitable distribution for $p_i$, hyperpriors for any parameters of this distribution and a prior for $N$. However two related things prevent us from doing this: (1) BUGS does not allow the index of the loop to be stochastic and (2) the dimension of $p$ changes each time we update $N$ to a new value.\n\n\\cite{Royle2007} provide a neat solution that is easily implemented in BUGS based on a superpopulation. However, a drawback of their model and code is that $N$ does not appear in the likelihood; rather it is a derived parameter. This makes specification of priors on $N$ more of a challenge than if $N$ appears explicitly in the likelihood. \n\nHere we contrast approaches to fitting $M_h$ using the models of \\cite{King2008} and \\cite{Royle2007} when using Bayesian methods. The approach taken by \\cite{King2008}, as well as \\cite{Durban2005,Schofield2008} and \\cite{Wright2009} for similar models, is to consider a trans-dimensional (TD) algorithm \\cite[]{Carlin1995,Green1995} such as the reversible jump MCMC (RJMCMC) algorithm \\cite[]{Green1995}. In contrast, \\cite{Royle2007} as well as \\cite{Royle2008,Link2010,Schofield2010} use a data augmentation (DA) approach \\cite[]{Tanner1987} to fit the model. \n\nOur purpose is to:\n\\begin{itemize}\n \\item[1.] Clarify the distinction between the different likelihoods used to fit $M_h$ models by MCMC.\n \\item[2.] Provide simple BUGS code for $M_h$ that explicitly includes $N$ as a parameter.\n \\item[3.] Consider the comment made by \\cite{Royle2007} that their approach is not a case of trans-dimensional MCMC and fundamentally differs from approaches that use RJMCMC.\n\\end{itemize}\n\n\n\\section{$M_h$ Likelihoods}\\label{sect:comparelike}\nLikelihoods for model $M_{h}$ include an integrated likelihood \\cite[]{Burnham1978,Otis1978}, a ``full'' likelihood \\cite[]{Pledger2000, King2008}, or a super-population likelihood \\cite[]{Royle2007}. The three most obvious differences between these approaches are (i) whether or not $N$ is explicitly included in the likelihood, (ii) the presence or absence of combinatorial terms involving $N$, and (iii) the presence or absence of a superpopulation (with corresponding value $M$) in the model. The likelihood of \\cite{Royle2007} does not include $N$, or combinatorial terms involving $N$, but does include $M$. The full likelihood approach of \\cite{Pledger2000} explicitly includes $N$ and associated combinatorial terms, does not include $M$ and treats the vector of capture probabilities $\\bm{p}$ as parameters. \\cite{Burnham1978} explicitly includes $N$, does not include $M$ and models $\\bm p$ as random effects; combinatorial terms appear once the capture probabilities are integrated out of the likelihood.\n\n\n\\begin{flushright}{\\bf Table 1 about here} \\end{flushright}\n\\subsection{Burnham's Integrated Likelihood}\nA complete data model can be written as $\\Pr(\\bm X = \\bm x | \\bm P = \\bm p)$; as shorthand we use $[\\bm x | \\bm p]$. Both $\\bm x$ and $\\bm p$ are unobserved, although we know values of some rows of $\\bm x$ through the observed data matrix $\\bm x^{obs}$. Following \\cite{Burnham1978} we can write the model as \n\\begin{eqnarray}\n [\\bm x | \\bm p] &=& \\prod_{i=1}^N \\prod_{j=1}^k p_i^{x_{ij}}(1-p_i)^{1-x_{ij}} \\label{eq:XgivenP}\n\n\\end{eqnarray}\nunder the assumption that, conditional on $\\bm P = \\bm p$, capture events are independent among individuals and occasions. Conditional on $N$, which is treated as a parameter, the number of rows of $\\bm x$ are now known. \n\nAs noted by \\cite{Burnham1978} this model is over-parameterized and not useful for estimation. The usual remedy for over-parameterization is to model $\\bm P$ as random effects sampled from a distribution with $pdf$ $f_{\\bm P}(\\bm p | \\theta_p)$ for $\\bm P \\in \\mathbb{R}_{[0,1]}^N$. We then integrate over $\\bm p$ to obtain the integrated likelihood given by \\cite{Burnham1978}; this is the observed data likelihood \\cite[]{Gelman2004}.\n\n\\subsection{Complete Data Likelihood (CDL)}\nThe model (\\ref{eq:XgivenP}) is underspecified.\nAdding the random effects distribution for $\\bm P$ we can write a CDL as\n\\begin{eqnarray}\n [\\bm x | \\bm p][\\bm p | \\theta_p] = [\\bm p |\\theta_p] \\times \\prod_{i=1}^N \\prod_{j=1}^k p_i^{x_{ij}}(1-p_i)^{1-x_{ij}} \\label{eq:CDL}.\n\\end{eqnarray}\nAs before, both $\\bm x$ and $\\bm p$ are unobserved. Associated with $\\bm x^{obs}$ is a subspace of ${\\mathcal S}$, which we denote by $\\Omega(\\bm x)$, of capture history matrices that are identical in appearance to $\\bm x^{obs}$. The only differences among elements of $\\Omega(\\bm x)$ are that the row indices are permuted. \n\nUnder the assumption that $x_i | p_i$ and $x_j | p_j$ are independent for all $i \\neq j$, it follows that the $x$'s are exchangeable given the $p$'s and that the elements in $\\Omega(\\bm x)$ have identical joint probabilities across their rows. A heuristic explanation is that once we have drawn $(x_i, p_i)$ pairs from their joint distribution they are linked -- when we permute the $x$'s the $p$'s get permuted with them.\n\nIt now follows that\n\\begin{eqnarray}\n [N, \\bm p | \\bm x^{obs}, \\theta_p] &\\propto& [\\bm x^{obs} | \\bm p][\\bm p | \\theta_p] \\nonumber \\\\\n &=&\\frac{N!}{\\prod_{h \\in {\\mathcal H}}z_h!} \\times \\prod_{i=1}^N [p_i | \\theta_p] \\prod_{j=1}^k p_i^{x_{ij}}(1-p_i)^{1-x_{ij}} \\label{eq:CDL2}\n\\end{eqnarray}\nis also a likelihood where $\\frac{N!}{\\prod_{h \\in {\\mathcal H}} z_h!}$ is the dimension of $\\Omega(\\bm x)$. The likelihood (\\ref{eq:CDL2}) is not the observed data likelihood (ODL) as it still contains unobserved elements. If we integrate across the random effects we obtain the ODL given by \\cite{Burnham1978}. However, the likelihood (\\ref{eq:CDL2}) is useful in Bayesian inference carried out by MCMC in that explicit modelling in terms of latent components can be a useful device for improving mixing.\n\n\n\\subsection{Super Population CDL}\n\\cite{Royle2007} describe a CDL in terms of a super-population $M$ (essentially an upper bound on the population size $N$), with $\\bm X$ now an $M$ by $k$ matrix. The idea is that we replace $N$ in the likelihood with the outcomes $\\bm w$ of the random vector $\\bm W$ of dimension $M$ where $W_{i}$ takes the value $1$ if individual $i$ is included in the population and $0$ otherwise such that $N=\\sum_{i=1}^{M}W_{i}$. By assumption, $\\bm W$ and $\\bm P$ are independent given their parameters. The complete model on which the \\cite{Royle2007} approach is based can be written as\n\\begin{eqnarray}\n [\\bm x, \\bm w, \\bm p | \\alpha_p,\\psi] &=& [\\bm w | \\psi][\\bm p | \\bm w, \\alpha_p][\\bm x |\\bm w, \\bm p] \\nonumber \\\\\n &=& [\\bm w | \\psi][\\bm p | \\alpha_p][\\bm x |\\bm w, \\bm p] \\nonumber\\\\\n &=&\\prod_{i=1}^M \\prod_{j=1}^t (p_iw_i)^{x_{ij}} (1-p_iw_i)^{1-x_{ij}} f_P(p_i) \\psi^{w_i}(1-\\psi)^{1-w_i} \\nonumber \\\\\n &=&\\prod_{i=1}^M (p_iw_i)^{y_i} (1-p_iw_i)^{t-y_i} f_P(p_i) \\psi^{w_i}(1-\\psi)^{1-w_i} \\label{eq:spCDL}.\n\\end{eqnarray}\nThe sample space ${\\mathcal S}$ for $\\bm x$ is now of dimension $2^{Mt}$ and the sample space ${\\mathcal W}$ for $\\bm w$ of dimension $2^M$. The sample space for $\\bm p$ is given by ${\\mathcal P} =\\mathbb{R}_{[0,1]}^M$. \n\nDenoting the $cdf$ for $p$ by $F(p) = F_{\\theta_p}(p)$ and integrating over $\\bm p$ and $\\bm w$, we obtain\n\\[\n [\\bm x^{obs}, N | \\alpha_p, \\psi, M] = \\frac{N!}{\\prod_{h \\neq \\bm 0}z_h!(N-n)!} \\prod_{j=0}^t \\theta_{jF}^{f_j} \\times \\binom{M}{N}\\psi^N(1-\\psi)^{M-N},\n\\]\nwhere $\\theta_{jF} = \\int_0^1 p^j(1-p)^{t-j}dF(p).$ That is, we have Burnham's marginal likelihood multiplied by the prior for $N$ induced by the model $W_i \\sim Bern(\\psi)$. \n\nInference in the two approaches will be identical for the appropriate choice of prior on $\\psi$ in the \\cite{Royle2007} likelihood. For example, \n\\[\n f(\\psi) \\propto 1\n\\]\ncorresponds to a discrete uniform prior on $N$ as shown by \\cite{Royle2007}. Another choice could be:\n\\[\n f(\\psi) \\propto \\frac{1}{\\psi}\n\\]\nwhich leads to $N\\propto \\frac{1}{N}$ but with the upper bound $M$. This is Jeffrey's prior for $N$ but truncated to the range $[0, M]$. In general, the prior on $N$ induced in the super-population approach is the marginal distribution for $N$ arising from a binomial mixture, with mixing distribution given by the prior on $\\psi$. While it is relatively easy to construct reference priors of the form given above, constructing general priors on $N$ that are not binomial mixtures will be more difficult. We do not wish to overstate this problem as the binomial mixtures provide a flexible class of priors including the beta-binomial and the negative-binomial as a limiting case. \n\n\\section{Gibbs Sampling and Model $M_h$}\n\\cite{Durban2005},\\cite{King2008} and \\cite{Schofield2008} all used a trans-dimensional (RJMCMC) approach to fit $M_h$ using the CDL described above. \\cite{Wright2009} used a similar approach in an application where individual identities were uncertain. In this approach, sampling from the full conditional for $N$ includes an explicit reversible-jump step in the sense of \\cite{Green1995}. In contrast, the reversible-jump step is unecessary in the super-population approach of \\cite{Royle2007}. \n\nTo describe what we call the ``standard approach'' to Gibbs sampling we start with the complete data likelihood (\\ref{eq:CDL}). Although we describe model fitting from a Bayesian stand-point, we could implement a Frequentist approach using an EM-algorithm instead of Gibbs sampling.\n\nFor simplicity we assume $\\theta_{p}$ is known. A Gibbs sampler can be constructed by alternating sampling from the following distributions,\n\\begin{enumerate}\n \\item The full conditional distribution for $p_{i}$, $\\forall i$:\n \\[\n [p_{i}|x,N,\\theta_{p}] \\propto p_{i}^{x_{ij}}(1-p_{i})^{1-x_{ij}}f(p|\\theta_{p}).\n \\]\n Depending on $[p|\\theta_{p}]$, we can either sample from this distribution directly, or use another a simulation based sampler, such as the Metropolis-Hastings algorithm or rejection sampling.\n \\item The full conditional distribution for $N$:\n \\[\n [N|x,p] \\propto\\frac{N!}{(N-n)!}\\prod_{i=1}^{N}\\prod_{j=1}^{k}p_{i}^{x_{ij}}(1-p_{i})^{1-x_{ij}}f(N).\n \\]\n This update is more difficult since the value of $N$ changes the dimension of $p$, so we use a TD algorithm. Following \\cite{Schofield2008} we use a special case of the RJMCMC algorithm. \n We first propose a new value $N^{\\star}$ from $J(N^{\\star}|N)$ and specify $p^{\\star}$ to be a vector of length $N^{\\star}$ corresponding to the capture probabilities for the $N^{\\star}$ individuals. We set $p^{\\star}_{i}=p_{i}$ for all $i = 1,\\ldots,N$ and generate $p^{\\star}_{i}$ for $i = N+1,\\ldots,N^{\\star}$ from the prior distribution $[p|\\theta_{p}]$. Generating $p^{\\star}$ in such a way simplifies the RJMCMC algorithm. The acceptance probability is\n \\[\n q = \\frac{N^{\\star}!(N-n)!}{(N^{\\star}-n)!N!}\\prod_{i=N+1}^{N^{\\star}}(1-p^{\\star}_{i})^{k}\\frac{J(N|N^{\\star})}{J(N^{\\star}|N)}\\frac{f(N^{\\star})}{f(N)}.\n \\]\n If our proposed value $N^{\\star}n$. Here we chose to use a discrete uniform prior on $N$ so the values \\texttt{pee[i]} are specified as data and are all set to $\\frac{1}{M}$. Note that (i) there are many different ways we can parameterize this distribution and that this may affect the mixing speed as well as the time taken to run the code.\n \\item[Line 13] Including the term $\\propto \\frac{N!}{(N-n)!}$ using the approximation of \\cite{Durban2005} who treat $n$ as binomially distributed with index $N$ and success probability $\\epsilon \\approx 0$. We assume that tags are correctly read and so do not include $\\prod y_{h}!$. Another approach to including the normalizing constant is to include a variable {\\tt zero=0} and then substitute line 13 with two lines: (i) {\\tt zeros\\~{}dpois(lam)} and (ii) {\\tt lam <- -(logfact(N)-logfact(N-n)-logfact(M))}.\n\\end{description}\nNote that this model statement is appropriate for both BUGS and JAGS. We provide alternative model statements and the data and initial values required for both BUGS and JAGS for each of these modeling statements at \\url{www.maramatanga.com}.\n\\caption{BUGS Code For Model Mh Using RJMCMC}\n\\label{fig:rjbug}\n\\end{figure}\n\n\n\\begin{figure}[htbp]\n\\begin{minipage}{\\textwidth}\n \\begin{center}\n \\includegraphics[angle=270,width = \\textwidth]{Nplot.eps} \n \\caption{Traceplot and Posterior Density for $N$ for the snowshoe hare data fitted in JAGS using the code in figure \\ref{fig:rjbug}.}\n \\label{fig:JAGSout}\n \\end{center}\n\\end{minipage}\n\\end{figure}\n\n\\end{document}","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction} \n\nThe polaron is a model of an electron interacting with its self-induced polarization field of the underlying crystal. The description of the polarization as a quantum field corresponds to the Fr\\\"ohlich model \\cite{frohlich}. In the classical approximation, on the other hand, the dynamics of a polaron is described by the Landau--Pekar (LP) equations. For $(\\psi_t, \\varphi_t) \\in H^1(\\mathbb{R}^{3})\\times L^2(\\mathbb{R}^{3})$, where $\\psi_t$ is the electron wave function and $\\varphi_t$ denotes the phonon field, these equations read in suitable units\n\\begin{equation}\\label{eq:LP}\n\\begin{aligned}\ni \\partial_t \\psi_t & = h_{ \\varphi_t } \\psi_t, \\\\\ni \\alpha^2 \\partial_t \\varphi_t & = \\varphi_t + \\sigma_{\\psi_t},\n\\end{aligned}\n\\end{equation}\nwhere $h_{\\varphi}$ is the Schr\\\"odinger operator \n\\begin{equation}\nh_{\\varphi} = - \\Delta + V_\\varphi\n\\end{equation}\nwith potential \n\\begin{equation}\nV_\\varphi(x) = 2 (2\\pi)^{3\/2} \\Re\\, [ ( - \\Delta)^{-1\/2} \\varphi ] (x)= 4 (2\\pi)^{5\/2} |x|^{-2} * \\Re\\, \\varphi ,\n\\end{equation}\nand\n\\begin{equation}\n\\sigma_{\\psi} (x) = (2\\pi)^{3\/2}\\left[ (- \\Delta)^{-1\/2} \\vert \\psi \\vert^2\\right] (x) = 2(2\\pi)^{5\/2} |x|^{-2} * |\\psi|^2 .\n\\end{equation}\nThe parameter $\\alpha>0$ quantifies the strength of the coupling of the electron's charge to the polarization field. \n\nThe LP equations can be derived from the dynamics generated by the (quantum) Fr\\\"ohlich Hamiltonian for suitable initial states in the strong coupling limit $\\alpha\\to \\infty$ \\cite{LMRSS} (see also \\cite{FS,FG,Griesemer,LRSS,M} for earlier results on this problem). One of the outstanding open problems concerns the polaron's effective mass \\cite{LS,seiringer2020polaron,spohn1987effective}: due to the interaction with the polarization field, the electron effectively becomes heavier and behaves like a particle with a larger mass. This mass increases with the coupling $\\alpha$, and is expected to diverge as $\\alpha^4$ as $\\alpha\\to\\infty$. A precise asymptotic formula was obtained by Landau and Pekar \\cite{ landau1948effective} based on the classical approximation, and hence it is natural to ask to what extent the derivation of the LP equations in \\cite{LMRSS} allows to draw conclusions on the effective mass problem.\n\nIt is, however, far from obvious how to rigorously obtain the effective mass even on the classical level, i.e., from the LP equations \\eqref{eq:LP}. A heuristic derivation, reviewed in Section \\ref{sec:TW} below, considers traveling wave solutions of \\eqref{eq:LP} for non-zero velocity $v\\in \\mathbb{R}^3$, and expands the corresponding energy for small $v$. The existence of such solutions remains unclear, however, and we in fact conjecture that no such solutions exist for non-zero $v$. This is related to the fact the energy functional corresponding to \\eqref{eq:LP} (given in Eq. \\eqref{def:G} below) does not dominate the total momentum, and a computation of the ground state energy as a function of the (conserved) total momentum would simply yield a constant function (corresponding to an infinite effective mass). Due to the vanishing of the sound velocity in the medium, a moving electron can be expected to be slowed down to zero speed by emitting radiation. (See \\cite{Fr1,Fr2} for the study of a similar effect in a model of a classical particle coupled to a field.) \n\nIn this paper, we provide a novel definition of the effective mass for the LP equations. We shall argue that all low energy states have a well-defined notion of (initial) velocity, and hence we can minimize the energy functional among states with given velocity. Expanding the resulting energy-velocity relation for small velocity gives a definition of the effective mass, which coincides with the prediction by Landau and Pekar \\cite{landau1948effective}. \n\n\\subsection{Structure of the paper}\n In Section \\ref{sec:result}, we explain our rigorous approach to derive the energy-velocity relation of the system, allowing for a precise definition and computation of the effective mass. After introducing some notation and recalling fundamental properties of the Pekar energy functional in Section \\ref{sec:prel}, we identify in Section \\ref{sec:vel} a set of initial data for the LP equations for which it is possible to define the position, and consequently the velocity, at any time. We then arrive at an energy-velocity relation by defining $E(v)$ in Section \\ref{sec:initial} as the minimal energy among all admissible initial states of fixed initial velocity $v$. Finally, in Section \\ref{sec:thm} we state our main result, an expansion of $E(v)$ for small velocities $v$, allowing for the computation the effective mass of the system. \n\nSection \\ref{sec:proof-thm} contains the proof of our main result, Theorem \\ref{thm:freephase}. \n\nIn Section \\ref{sec:approaches} we discuss the formal approach to the effective mass via traveling waves. Moreover, we investigate an alternative definition of the effective mass, through an alternative notion of velocity of low-energy states.\n\n\n\n\n\\iffalse\n\nOur main object of study is the effective mass of the polaron in this classical approximation. Several computations are present in the physics literature concerning this quantity \\cite{alexandrov2010advances,landau1948effective,spohn1987effective}, which, however, are all based on heuristic arguments. In particular, it is conjectured that the effective mass diverges as $\\alpha^4$ as $\\alpha \\to \\infty$, with an explicit prefactor related to the Pekar minimizer $\\varphi_{\\rm P}$, i.e.\n\\begin{align}\n\\label{eq:mass-pred}\n\\lim_{\\alpha \\rightarrow \\infty} \\alpha^{-4} \\; m_{\\text{eff}} = \\frac{2}{3}\\| \\nabla \\varphi_{\\rm P} \\|^2 \\; .\n\\end{align}\nThis work provides a rigorous definition and computation of the effective mass, verifying \\eqref{eq:mass-pred}. Our approach relies on the introduction of a well-defined notion of position and velocity of the system and a perturbative expansion of the energy functional associated to the Landau-Pekar equations. These allow to derive an explicit energy-velocity relation for the system, yielding the exact value of the effective mass of the system. \n\n\\fi\n\n\n\\section{Main Results}\n\\label{sec:result}\n\n\\subsection{Preliminaries}\n\\label{sec:prel} We start by introducing further notation and recalling some known results. \nThe classical energy functional corresponding to the Landau--Pekar equations \\eqref{eq:LP} is defined on $H^1(\\mathbb{R}^{3})\\times L^2(\\mathbb{R}^{3})$ as \n\\begin{equation}\\label{def:G}\n\\mathcal{G} ( \\psi, \\varphi ) = \\langle \\psi, h_\\varphi \\psi \\rangle + \\| \\varphi \\|_2^2 \\quad \\text{for} \\quad \\|\\psi\\|_2=1 .\n\\end{equation}\nEquipped with the symplectic form $\\frac 1{2i} \\int d\\psi \\wedge d\\bar\\psi + \\frac{\\alpha^2}{2i} \\int d\\varphi\\wedge d\\bar\\varphi$, it defines a dynamical system leading to the LP equations \\eqref{eq:LP}. Moreover, $\\mathcal{G}$ is conserved along solutions of \\eqref{eq:LP}.\n\nIt was proved in \\cite{lieb1977existence} that the Pekar ground state energy \n\\begin{align}\n\\min \\mathcal{G}(\\psi,\\varphi) =: e_{\\rm P} \n\\end{align}\nis attained for the Pekar minimizers $( \\psi_{\\rm P}, \\varphi_{\\rm P})$, which are radial functions in $C^{\\infty}(\\mathbb{R}^{3})$ satisfying $\\psi_{\\rm P}>0$, $\\varphi_{\\rm P}=-\\sigma_{\\psi_{\\rm P}}$ and $\\psi_{\\rm P}=\\psi_{\\varphi_{\\rm P}}$, where $\\psi_{\\varphi}$ denotes the ground state of $h_{\\varphi}$ whenever it exists. Moreover, this minimizer is unique up to the symmetries of the problem, i.e., translation-invariance and multiplication of $\\psi$ by a phase. \nWe shall denote \n\\begin{align}\n\\label{eq:def-HP}\nH_{\\rm P} = h_{\\varphi_{\\rm P}} - \\mu_{\\rm P} \\quad \\text{with} \\quad \\mu_{\\rm P}=\\infspec h_{\\varphi_{\\rm P}} .\n\\end{align} \n\nAssociated to $\\mathcal{G}$, there are the two functionals \n\\begin{align}\n\\label{eq:EandF}\n\\mathcal{E}(\\psi):=\\inf_{\\varphi\\in L^2(\\mathbb{R}^{3})} \\mathcal{G}(\\psi,\\varphi), \\quad \\mathcal{F}(\\varphi):=\\inf_{\\psi\\in H^1(\\mathbb{R}^{3}) \\atop \\|\\psi\\|_2=1} \\mathcal{G}(\\psi,\\varphi)\n\\end{align}\nand clearly $e_{\\rm P}=\\min \\mathcal{G}(\\psi,\\varphi)=\\min\\mathcal{E}(\\psi)=\\min \\mathcal{F}(\\varphi)$. \nWe also define the manifolds of minimizers \n\\begin{align}\n\\label{eq:PekEn}\n\\mathcal{M}_{\\mathcal{G}}:=\\{(\\psi,\\varphi) \\,|\\, \\mathcal{G}(\\psi,\\varphi)&=e_{\\rm P}\\}, \\ \\mathcal{M}_{\\mathcal{E}}:=\\left\\{\\psi \\,|\\, \\mathcal{E}(\\psi)=e_{\\rm P}\\right\\}, \\ \\mathcal{M}_{\\mathcal{F}}:=\\left\\{\\varphi \\,|\\, \\mathcal{F}(\\varphi)=e_{\\rm P}\\right\\}.\n\\end{align}\nThe results in \\cite{lieb1977existence} imply that we can write these in terms of the Pekar minimizers $(\\psi_{\\rm P}, \\varphi_{\\rm P})$ as \n\\begin{align}\n\\mathcal{M}_{\\mathcal{G}}=\\{(e^{i\\theta}& \\psi_{\\rm P}^y,\\varphi_{\\rm P}^y) \\,|\\, \\theta \\in [0,2\\pi), \\, y \\in \\mathbb{R}^{3} \\},\\nonumber\\\\\n\\mathcal{M}_{\\mathcal{E}}=\\{&e^{i\\theta} \\psi_{\\rm P}^y \\,|\\, \\theta \\in [0,2\\pi), \\, y \\in \\mathbb{R}^{3}\\},\\nonumber\\\\\n&\\mathcal{M}_{\\mathcal{F}}=\\{\\varphi_{\\rm P}^y \\,|\\, y \\in \\mathbb{R}^{3}\\}\n\\end{align}\nwhere $f^y:=f(\\cdot-y)$ for any function $f$. Furthermore, it can be deduced from the results in \\cite{lenzmann2009uniqueness} that the energy functionals $\\mathcal{F}$ and $\\mathcal{E}$ are both coercive (see \\cite[Lemmas~2.6 and 2.7]{feliciangeli2021persistence}), i.e., there exists $C>0$ such that\n\\begin{align}\n\\label{eq:coercivity}\n\\mathcal{F}(\\varphi)\\geq e_{\\rm P}+C\\dist^2_{L^2}(\\varphi,\\mathcal{M}_{\\mathcal{F}}), \\quad \\mathcal{E}(\\psi)\\geq e_{\\rm P}+ C \\dist^2_{H^1}(\\psi,\\mathcal{M}_{\\mathcal{E}}).\n\\end{align} \n\n\nThe following Lemma on properties of the projection onto the manifold $\\mathcal{M}_{\\mathcal{F}}$ will be important for our analysis below. Its proof will be given in Appendix \\ref{app:projection}.\n\n\\begin{lem}\n\t\\label{lem:minEFproj}\n\tThere exists $\\delta>0$ such that the $L^2$-projection onto $\\mathcal{M}_{\\mathcal{F}}$, is well-defined (i.e., unique) on \n\t\\begin{align}\n\t(\\mathcal{M}_{\\mathcal{F}})_{\\delta}&:=\\{\\varphi\\in L^2(\\mathbb{R}^{3})\\;|\\; \\dist_{L^2}(\\varphi, \\mathcal{M}_{\\mathcal{F}})\\leq \\delta\\} \\; .\n\t\\end{align}\n\tFor any $\\varphi \\in (\\mathcal{M}_{\\mathcal{F}})_{\\delta}$, we define $z_\\varphi \\in \\mathbb{R}^3$ via \n\t\\begin{align}\n\tP_{L^2}^{\\mathcal{M}_{\\mathcal{F}}}(\\varphi)&=\\varphi_{\\rm P}^{z_{\\varphi}} \\; .\n\t\\end{align}\n\tThen $z_{\\varphi}$ is a differentiable function from $(\\mathcal{M}_{\\mathcal{F}})_{\\delta}$ to $\\mathbb{R}^{3}$ and its partial derivative in the direction $\\eta \\in L^2(\\mathbb{R}^3)$ is given by \n\t\\begin{align}\n\t \\partial_t z_{\\varphi + t \\eta} \\restriction_{t=0}\n\t = A_\\varphi^{-1} \\langle \\Re \\eta \\vert \\nabla \\varphi_P^{z_\\varphi} \\rangle,\n\t\\end{align} \n\twhere $A$ is the invertible matrix defined for any $\\varphi \\in(\\mathcal{M}_{\\mathcal{F}})_{\\delta}$ by $A_{i,j} := - \\Re \\langle \\varphi \\vert \\partial_i \\partial_j \\varphi_P^{z_\\varphi} \\rangle$. \n\\end{lem} \n\n\n\\begin{rem} \\label{rem:E} Likewise, it can be shown that the $H^1$- (resp. $L^2$-) projection onto $\\mathcal{M}_{\\mathcal{E}}$ have similar properties: There exists $\\delta >0 $ such that the $H^1$- (resp. the $L^2$-) projection onto $\\mathcal{M}_{\\mathcal{E}}$ \n\t\\begin{align}\n\t\\label{eq:proj-E}\n\tP_{H^1}^{\\mathcal{M}_{\\mathcal{E}}}(\\psi)=e^{i\\theta_{\\psi}}\\psi_{\\rm P}^{y_{\\psi}}, \\quad\\Big(\\text{resp.} \\ P_{L^2}^{\\mathcal{M}_{\\mathcal{E}}}(\\psi)&=e^{i\\theta'_{\\psi}}\\psi_{\\rm P}^{y'_{\\psi}} \\Big) \n\t\\end{align}\n\tis well-defined on the set $(\\mathcal{M}_{\\mathcal{E}})^{H^1}_{\\delta}:=\\{\\psi\\in L^2(\\mathbb{R}^{3})\\;|\\; \\dist_{H^1}(\\psi, \\mathcal{M}_{\\mathcal{E}})\\leq \\delta\\}$ (resp. $(\\mathcal{M}_{\\mathcal{E}})^{L^2}_{\\delta}:=\\{\\psi\\in L^2(\\mathbb{R}^{3})\\;|\\; \\dist_{L^2}(\\psi, \\mathcal{M}_{\\mathcal{E}})\\leq \\delta\\}$) and the functions $y_\\psi, \\theta_\\psi$ (resp. $y'_\\psi, \\theta'_\\psi$) defined through \\eqref{eq:proj-E} \nare differentiable functions from $(\\mathcal{M}_{\\mathcal{E}})^{H^1}_{\\delta}$ (resp. $(\\mathcal{M}_{\\mathcal{E}})^{L^2}_{\\delta}$) to $\\mathbb{R} \/ (2\\pi \\mathbb{Z})$ and $\\mathbb{R}^{3}$. \n\\end{rem}\n\n\\subsection{Position and velocity of solutions} \\label{sec:vel}\n\nIn this section, we give a meaning to the notion of position, and therefore velocity, for solutions of the Landau--Pekar equations (at least for a class of initial data). There is a natural way of defining, given $\\psi_t$, the position of the electron at time $t$, which is simply given by \n\\begin{align}\n\\label{eq:defXel}\n\\text{X}_{\\text{el}}(t):=\\expval{x}{\\psi_t}.\n\\end{align}\nThis yields, by straightforward computations using \\eqref{eq:LP}, that\n\\begin{align}\\label{vel}\n\\text{V}_{\\text{el}}(t):=\\frac{d}{dt}\\text{X}_{\\text{el}}(t)=2\\expval{-i\\nabla}{\\psi_t}.\n\\end{align} \nNote that \\eqref{vel} is always well-defined for $\\psi\\in H^1(\\mathbb{R}^3)$, even although \\eqref{eq:defXel} not necessarily is. \n\nFor the phonon field, the situation is more complicated as $\\varphi$ cannot be interpreted as a probability distribution over positions. This calls for a different approach. By \\eqref{eq:coercivity}, Lemma \\ref{lem:minEFproj} and the conservation of $\\mathcal{G}$ along solutions of \\eqref{eq:LP}, we know that there exists $\\delta^*$ such that for any initial condition $(\\psi_0,\\varphi_0)$ such that \n\\begin{align}\n\\label{eq:IDuniqueproj}\n\\mathcal{G}(\\psi_0,\\varphi_0)\\leq e_{\\rm P}+\\delta^*,\n\\end{align}\n$\\varphi_t$\nadmits a unique $L^2$-projection $\\varphi_{\\rm P}^{z(t)}$ onto $\\mathcal{M}_{\\mathcal{F}}$ for all times.\nWe use this to define \n\\begin{align}\n\\label{eq:defXph}\n\\text{X}_{\\text{ph}}(t):=z(t), \\quad \\text{V}_{\\text{ph}}:=\\frac{d}{dt} \\text{X}_{\\text{ph}}(t)=\\dot{z}(t).\n\\end{align}\nNote that $\\text{X}_{\\text{ph}}(t)$ is indeed differentiable by Lemma \\ref{lem:minEFproj} and the differentiability of the LP dynamics. At this point, for any initial data satisfying \\eqref{eq:IDuniqueproj}, we have a well-defined notion of position and velocity for all times, admittedly in a much less explicit form for the phonon field.\n\n\\subsection{Initial conditions of velocity $v$}\n\\label{sec:initial}\n \nFor any $v\\in \\mathbb{R}^{3}$ (or at least for $|v|$ sufficiently small), we are now interested in considering all initial conditions $(\\psi_0,\\varphi_0)$ whose solutions have instantaneous velocity $v$ at $t=0$ (both in the electron and in the phonon coordinate) and to then minimize the functional $\\mathcal{G}$ over such states. This will give us an explicit relation between the energy and the velocity of the system, allowing us to define the effective mass of the polaron in the classical setting defined by the Landau--Pekar equations.\n\nNote that by radial symmetry of the problem only the absolute value of the velocity, and not its direction, affects our analysis. Hence, for $v\\in \\mathbb{R}$, we consider initial conditions $(\\psi_0,\\varphi_0)$ such that\n\\begin{itemize}\n\t\\item[(i)] $(\\psi_0, \\varphi_0) \\in H^1 (\\mathbb{R}^3) \\times L^2 ( \\mathbb{R}^3) $ with $\\| \\psi_0\\|_2=1$ and such that \\eqref{eq:IDuniqueproj} is satisfied,\n\t\\item[(ii)] $\\text{V}_{\\text{el}}(0)=\\text{V}_{\\text{ph}}(0)=v (1,0,0)$. \n\\end{itemize}\nThe set of admissible initial conditions of velocity $v\\in \\mathbb{R}$ can hence be compactly written as\n\\begin{align}\n\\label{eq:defI}\nI_v:= \\lbrace (\\psi_0, \\varphi_0) \\, \\vert \\, \\text{(i), (ii) are satisfied} \\rbrace.\n\\end{align}\nWe will show below that it is non-empty for small enough $v$. \n\n\\subsection{Expansion of the energy} \\label{sec:thm} In order to compute the effective mass of the polaron, we now minimize the energy $\\mathcal{G}$ over the set $I_v$.\nTo this end, we define the energy\n\\begin{align}\\label{def:Ev}\nE(v) :=\\inf_{(\\psi_0,\\varphi_0)\\in I_v} \\mathcal{G}(\\psi_0,\\varphi_0).\n\\end{align}\nThe following theorem gives an expansion of $E(v)$ for sufficiently small velocities $v$. Its proof will be given in Section \\ref{sec:proof-thm}.\n\n\\begin{theo}\n\t\\label{thm:freephase}\n\tAs $v \\rightarrow 0$ we have \n\t\\begin{equation}\n\t\\label{eq:thm}\n\tE(v)=e_{\\rm P}+v^2 \\left( \\frac{1}{4}+ \\frac{\\alpha^4}{3} \\| \\nabla \\varphi_{\\rm P} \\|_2^2 \\right) + O(v^3).\t\n\t\\end{equation}\n\\end{theo}\n\nSince the kinetic energy of a particle of mass $m$ and velocity $v$ equals $m v^2\/2$, \\eqref{eq:thm} identifies the effective mass of the system as\n\\begin{align}\n\\label{eq:OurEffMass}\nm_{\\text{eff}}= \\lim_{v \\rightarrow 0} \\frac{E(v) - e_{\\rm P}}{v^2\/2} = \\frac 1 2 +\\frac {2\\alpha^4} 3 \\| \\nabla \\varphi_{\\rm P} \\|_2^2 .\n\\end{align}\nThe first term $1\/2$ is simply the bare mass of the election, while the second term $\\frac {2\\alpha^4} 3 \\| \\nabla \\varphi_{\\rm P} \\|_2^2$ corresponds to the additional mass acquired through the interaction with the phonon field, and agrees with the prediction in \\cite{landau1948effective}.\n\n\n\\begin{rem}[Traveling waves]\n\\label{rem:TW}\nThe heuristic computations contained in the physics literature concerning $m_{\\text{eff}}$ \\cite{alexandrov2010advances,landau1948effective}\nall rely, in one way or another, on the existence of traveling wave solutions of the LP equations of velocity $v$ (at least for sufficiently small velocity), i.e. solutions with initial data $(\\psi_v,\\varphi_v)$ such that \n\\begin{align}\n(\\psi_t(x),\\varphi_t(x))=(e^{-ie_{v}t}\\psi_v(x-vt),\\varphi_v(x-vt))\n\\end{align}\nfor suitable $e_v \\in \\mathbb{R}$. \nSuch solutions would allow to define the energy of the system at velocity $v$ as $E^{\\text{TW}}(v) = \\mathcal{G}( \\psi_v, \\varphi_v)$, and a perturbative calculation (discussed in Section \\ref{sec:TW} below) yields indeed\n\\begin{align}\n\\label{eq:mass-predTW}\n\\lim_{v \\rightarrow 0} \\frac{E^{\\text{TW}}(v) - e_{\\rm P}}{v^2\/2} = \\frac 1 2 +\\frac {2\\alpha^4} 3 \\| \\nabla \\varphi_{\\rm P} \\|^2, \n\\end{align} \nin agreement with \\eqref{eq:OurEffMass}. \n Unfortunately, this approach turns out to be only formal, and we conjecture traveling wave solutions to not exist for any $\\alpha>0$, $v>0$.\n\\end{rem} \n\n\n\\iffalse\n\\begin{rem}[The Fr\\\"ohlich polaron model]\n\t\nWe remark that the true dynamics of the polaron is given by its quantum mechanical description based on the Fr\\\"ohlich Hamiltonian, denoted in the following by $H^F_{\\alpha}$, which models the polarization field as quantized phonons (we refer to \\cite{seiringer2020polaron} and references therein for the definition of $H^F_{\\alpha}$ and a comprehensive discussion of its properties and the model). It features a dimensionless coupling constant $\\alpha$, modeling the strength of the interaction between the electron and the crystal. In the strong coupling regime $\\alpha\\gg 1$, the quantum field displays classical behavior. In particular, it was shown (see \\cite{FS,FG,Griesemer,LMRSS,LRSS,M}) that the LP equations provide an effective description of its dynamics for large $\\alpha$. Therefore, we keep this model in mind as a motivation for our work, even if the analysis that we carry out on the classical level says little about the full quantum problem. \n\nNote that, in this context, very little is known about the effective mass of the system $m^F_{\\text{eff}}$. One approach to define it, relies on deriving a momentum-energy relation of the system (rather than the velocity-energy relation we opt for at the classical level). This approach exploits the translational invariance of the Hamiltonian, which allows to write it in its momentum fiber-integral decomposition as $H^F_\\alpha = \\int_{\\mathbb{R}^3}^\\oplus H^P_\\alpha \\; dP$, arriving at the relation $E(P) = \\inf \\text{spec}H^P_\\alpha $, which implicitly defines $m_{\\text{eff}}^{F}$. The only rigorous result concerning this definition is contained in \\cite{LS}, where it is shown that $m_{\\text{eff}}^{F}\\to \\infty$ as $\\alpha\\to \\infty$. \n\\end{rem}\n\\fi\n\n\n\n\\begin{rem} \\label{rem:AlternatePos}\nIn Section \\ref{sec:vel}, we used the standard approach from quantum mechanics to define the electron's position \\eqref{eq:defXel} and velocity \\eqref{vel}. We could, instead, use also for the electron a similar approach to the one we use for the phonon field (i.e. \\eqref{eq:defXph}) through the projection onto the manifold of minimizers $\\mathcal{M}_{\\mathcal{E}}$. A natural question is whether one obtains the same effective mass using this different notion of position. In Section \\ref{sec:diff-vel}, we show that, in fact, this alternate definition yields a different effective mass equal to \n\\begin{align}\n\\widetilde{m}_{\\text{eff}} = \\frac{2 \\| \\nabla \\psi_{\\rm P} \\|_2^4}{3 \\|\\nabla \\varphi_{\\rm P} \\|_2^2 } +\\frac{2\\alpha^4}{3} \\| \\nabla \\varphi_{\\rm P} \\|_2^2 .\n\\end{align}\nThis coincides with \\eqref{eq:OurEffMass} and \\eqref{eq:mass-predTW} for large $\\alpha$ (hence still confirming the prediction in \\cite{landau1948effective}), but differs in the $O(1)$ term. In fact, as we discuss in Section \\ref{sec:diff-vel}, one has $\\widetilde{m}_{\\text{eff}} < m_{\\text{eff}}$. \n\\end{rem}\n\n\\section{Proof of Theorem~\\ref{thm:freephase}}\\label{sec:proof-thm}\n\n\n\n Let us denote $\\delta_1=\\psi_0-\\psi_{\\rm P}$ and $\\delta_2=\\varphi_0-\\varphi_{\\rm P}$. Expanding $\\mathcal{G}$ in \\eqref{def:G} and using that $\\varphi_{\\rm P}=-\\sigma_{\\psi_{\\rm P}}$ we find\n\\begin{align}\n\\mathcal{G}(\\psi_0,\\varphi_0) & =\\mathcal{G}(\\psi_{\\rm P}+\\delta_1,\\varphi_{\\rm P}+\\delta_2) \\nonumber\\\\\n& =e_{\\rm P}+2\\bra{\\psi_{\\rm P}}h_{\\varphi_{\\rm P}}\\ket{\\Re \\delta_1}\\nonumber\\\\\n&\\quad +\\expval{h_{\\varphi_{\\rm P}}}{\\delta_1}+2\\bra{\\Re \\delta_1}V_{\\delta_2}\\ket{\\psi_{\\rm P}}+\\|\\delta_2\\|_2^2 +\\expval{V_{\\delta_2}}{\\delta_1}. \\label{eq:expansion-G}\n\\end{align}\nSince $\\psi_0$ is normalized, we have\n\\begin{align}\\label{pdd}\n1=\\|\\psi_0\\|_2^2=\\|\\psi_{\\rm P}+\\delta_1\\|_2^2=1+\\|\\delta_1\\|_2^2+2\\bra{\\psi_{\\rm P}}\\ket{\\Re \\delta_1}\\iff 2\\bra{\\psi_{\\rm P}}\\ket{\\Re \\delta_1}=-\\|\\delta_1\\|_2^2.\n\\end{align}\nHence\n\\begin{align}\n2\\bra{\\psi_{\\rm P}}h_{\\varphi_{\\rm P}}\\ket{\\Re \\delta_1}=2\\mu_{\\rm P}\\bra{\\psi_{\\rm P}}\\ket{\\Re \\delta_1}=-\\mu_{\\rm P}\\|\\delta_1\\|_2^2,\n\\end{align}\nand using $\\| V_{\\delta_2} \\delta_1 \\|_2 \\leq C \\| \\delta_2\\|_2 \\; \\|\\delta_1 \\|_{H^1}$ (see, e.g., \\cite[Lemma III.2]{LRSS}) we arrive at\n\\begin{align}\n\\mathcal{G}(\\psi_0,\\varphi_0)=e_{\\rm P}+\\expval{H_{\\rm P}}{\\delta_1}+2\\bra{\\Re \\delta_1}V_{\\delta_2}\\ket{\\psi_{\\rm P}}+\\|\\delta_2\\|_2^2+O(\\|\\delta_2\\|_2 \\|\\delta_1\\|_{H^1}^2).\n\\end{align}\nBy completing the square, we have\n\\begin{align}\n&\\|\\Re\\delta_2\\|_2^2+2\\bra{\\Re \\delta_1}V_{\\delta_2}\\ket{\\psi_{\\rm P}}\\nonumber\\\\\n&=\\|\\Re \\delta_2 +2 (2\\pi)^{3\/2}(-\\Delta)^{-1\/2} (\\psi_{\\rm P}\\Re \\delta_1)\\|_2^2-4 (2\\pi)^3\\expval{\\psi_{\\rm P}(-\\Delta)^{-1} \\psi_{\\rm P}}{\\Re \\delta_1}\n\\end{align}\nand therefore\n\\begin{align}\\label{36}\n\\mathcal{G}(\\psi_0,\\varphi_0) &= e_{\\rm P}+\\expval{H_{\\rm P}}{\\Im\\psi_0}+\\|\\Im \\varphi_0\\|_2^2 \\nonumber \\\\ & \\quad +\\|\\Re \\delta_2 +2(2\\pi)^{3\/2}(-\\Delta)^{-1\/2} (\\psi_{\\rm P}\\Re \\delta_1)\\|_2^2\\nonumber\\\\\n&\\quad +\\expval{H_{\\rm P}-4X_{\\rm P}}{\\Re \\delta_1}+O(\\|\\delta_2\\|_2 \\|\\delta_1\\|^2_{H^1}) ,\n\\end{align} \nwhere $X_{\\rm P}$ is the operator with integral kernel $X_{\\rm P}(x,y):= (2\\pi)^3\\psi_{\\rm P}(x)(-\\Delta)^{-1}(x,y) \\psi_{\\rm P}(y)$. Since $X_{\\rm P}$ is bounded, and $\\| P_{\\psi_{\\rm P}} \\Re\\delta_1\\| = \\|\\delta_1\\|_2^2\/2$ by \\eqref{pdd} (with $P_{\\psi_{\\rm P}} = |\\psi_{\\rm P}\\rangle\\langle \\psi_{\\rm P}|$ the rank one projection onto $\\psi_{\\rm P}$), \nwe also have \n\\begin{align}\n\\mathcal{G}(\\psi_0,\\varphi_0)& =e_{\\rm P}+\\expval{H_{\\rm P}}{\\Im\\psi_0}+\\|\\Im\\varphi_0\\|_2^2 \\nonumber \\\\ & \\quad +\\|\\Re \\delta_2 +2(2\\pi)^{3\/2} (-\\Delta)^{-1\/2} (\\psi_{\\rm P}\\Re \\delta_1)\\|_2^2\\nonumber\\\\\n &\\quad+\\expval{Q(H_{\\rm P}-4X)Q}{\\Re \\delta_1}+O(\\|\\delta_2\\|_2\\|\\delta_1\\|^2_{H^1})+O(\\|\\delta_1\\|_{L^2}^3) , \\label{eq:expressionG}\n\\end{align}\nwhere $Q=1\\!\\!1-P_{\\psi_{\\rm P}}$. \n\n\n\n\\textbf{Upper Bound:} For sufficiently small $v$, we use as a trial state\n\\begin{align}\n\\label{eq:trialstate}\n(\\bar\\psi_0,\\bar\\varphi_0)=\\left( f_v \\psi_{\\rm P}+i g_v H_{\\rm P}^{-1}\\partial_1 \\psi_{\\rm P} , \\; \\varphi_{\\rm P}+iv\\alpha^2 \\partial_1 \\varphi_{\\rm P}\\right)\n\\end{align}\nwith $f_v, g_v >0$ given by \n\\begin{align}\n\\label{eq:defFG}\nf_v^2 := \\frac{2 v^2 \\|H_{\\rm P}^{-1}\\partial_1 \\psi_{\\rm P}\\|_2^2 }{1 - \\sqrt{1-4v^2 \\|H_{\\rm P}^{-1}\\partial_1 \\psi_{\\rm P}\\|_2^2 } }, \\; g_v^2 := \\frac{ 1 - \\sqrt{1-4v^2 \\|H_{\\rm P}^{-1}\\partial_1 \\psi_{\\rm P}\\|_2^2 } }{2 \\|H_{\\rm P}^{-1}\\partial_1 \\psi_{\\rm P}\\|_2^2 } .\n\\end{align}\nNote that $\\partial_1 \\psi_{\\rm P}$ is orthogonal to $\\psi_{\\rm P}$, hence $H_{\\rm P}^{-1} \\partial_1 \\psi_{\\rm P}$ is well-defined. \nWe begin by showing that \\eqref{eq:trialstate} is an element of the set of admissible initial data $I_v$ in \\eqref{eq:defI}. To prove that $(\\bar\\psi_0,\\bar\\varphi_0)$ satisfies (i), we only need to check that $\\bar\\psi_0$ is normalized (which follows easily from \\eqref{eq:defFG}) as the condition \\eqref{eq:IDuniqueproj} will follow a posteriori from the energy bound we shall derive. We now proceed to show that $(\\bar\\psi_0,\\bar\\varphi_0)$ satisfies (ii). For the electron, using that $H_{\\rm P}^{-1} \\partial_j \\psi_{\\rm P}= -x_j \\psi_{\\rm P} \/2 $ (which can be checked by applying $H_P$ and using that $[H_P,x_1]=-2 \\partial_1$) and consequently that \n\\begin{align}\n\\label{eq:sandwich}\n\\bra{\\partial_i \\psi_{\\rm P}}H_{\\rm P}^{-1}\\ket{\\partial_j \\psi_{\\rm P}}=\\delta_{ij}\/ 4 \n\\end{align} \nsince $\\psi_{\\rm P}$ is radial, \nwe can conclude that\n\\begin{align}\n\\label{eq:trial-v}\n-2\\expval{i \\partial_j }{\\bar \\psi_0}=4 f_v g_v \\langle H_{\\rm P}^{-1} \\partial_1 \\psi_{\\rm P} \\vert \\partial_j \\psi_{\\rm P} \\rangle =v \\delta_{j1} ,\n\\end{align}\ni.e., that $\\text{V}_{\\text{el}}(0)=v(1,0,0) $, as required.\n\nFor the phonons, we first note that $\\text{X}_{\\text{ph}}(0)=0$, since $\\Re \\bar{\\varphi}_0 = \\varphi_{\\rm P}$. Next, we derive a relation for the velocity of the phonons $\\text{V}_{\\text{ph}}(t) = \\dot{z}(t)$ in terms of their position $\\text{X}_{\\text{ph}} (t) =z(t)$ for general time $t$. Since \n\\begin{align}\n\\label{eq:defZ}\n\\min_z\\|\\varphi_t -\\varphi_{\\rm P}^z\\|_2^2=\\|\\varphi_t-\\varphi_{\\rm P}^{z(t)}\\|_2^2,\n\\end{align}\nthe position $z(t)$ necessarily has to satisfy \n\\begin{align}\n\\label{eq:orthPhon}\n\\Re \\bra{\\varphi_t}\\ket{(u\\cdot \\nabla) \\varphi_{\\rm P}^{z(t)}}=0 \\quad \\text{for all} \\quad u \\in \\mathbb{S}^2 \\iff \\Re \\varphi_t \\perp \\spn\\{\\nabla \\varphi_{\\rm P}^{z(t)}\\}.\n\\end{align}\nDifferentiating \\eqref{eq:orthPhon} w.r.t. time, using \\eqref{eq:LP} and evaluating the resulting expression at $t=0$, we arrive at \n\\begin{align}\n\\label{eq:defVph2}\n0=&\\Re\\bra{-i\\alpha^{-2} (u\\cdot \\nabla)(\\bar\\varphi_0+\\sigma_{\\bar\\psi_0})}\\ket{\\varphi_{\\rm P}}-\\Re\\bra{(\\dot{z}(0)\\cdot\\nabla)\\bar\\varphi_0}\\ket{(u\\cdot \\nabla)\\varphi_{\\rm P}}\\nonumber\\\\\n=&\\bra{-\\alpha^{-2} \\Im\\bar\\varphi_0}\\ket{(u\\cdot \\nabla)\\varphi_{\\rm P}}-\\bra{(\\dot{z}(0)\\cdot\\nabla)\\Re \\bar\\varphi_0}\\ket{(u\\cdot \\nabla)\\varphi_{\\rm P}}\\nonumber\\\\\n=&-\\bra{v\\partial_1 \\varphi_{\\rm P}}\\ket{(u\\cdot \\nabla)\\varphi_{\\rm P}}-\\bra{(\\dot{z}(0)\\cdot\\nabla)\\varphi_{\\rm P}}\\ket{(u\\cdot \\nabla)\\varphi_{\\rm P}},\n\\end{align}\nwhich the velocity $\\dot{z} (0)$ has to satisfy for all $u\\in \\mathbb{S}^2$, given its position $\\text{X}_{\\text{ph}}(0)=z (0)=0$. By invertibility of the coefficient matrix, \\eqref{eq:defVph2} has the unique solution $\\dot{z}(0)=v(1,0,0)$, and we indeed conclude that $\\text{V}_{\\text{ph}} (0) = v (1,0,0)$.\n\nWe now evaluate $\\mathcal{G}(\\bar\\psi_0,\\bar\\varphi_0)$. Since $f_v = 1 + O(v^2), g_v = v+ O(v^3)$, using \\eqref{eq:expressionG} and \\eqref{eq:sandwich} we find\n\\begin{align}\nE(v)\\leq \\mathcal{G}(\\bar\\psi_0,\\bar\\varphi_0)=e_{\\rm P}+v^2 \\left( \\frac 1 4+\\alpha^4 \\| \\partial_1 \\varphi_{\\rm P} \\|_2^2 \\right) + O(v^3)\n\\end{align}\nverifying on the one hand \\eqref{eq:IDuniqueproj} for sufficiently small $v$, and on the other hand the r.h.s. of \\eqref{eq:thm} as an upper bound on $E(v)$ (using that $\\varphi_{\\rm P}$ is radial).\n\n\n\n\\textbf{Lower Bound:} We first observe that to derive a lower bound on $E(v)$ we can w.l.o.g. restrict to initial conditions $(\\psi_0,\\varphi_0)$ satisfying additionally\n\\begin{align}\n\\label{eq:centeringI}\nP_{L^2}^{\\mathcal{M}_{\\mathcal{E}}}(\\psi_0)>0, \\\\\n\\label{eq:centeringII} \\text{X}_{\\text{ph}}(0)=0.\n\\end{align} \nThis simply follows from the invariance of $\\mathcal{G}$ under translations (of both $\\psi$ and $\\varphi$) and under changes of phase of $\\psi$. \nMoreover, by the upper bound derived in the first step of this proof and the coercivity of $\\mathcal{E}$ and $\\mathcal{F}$ in \\eqref{eq:coercivity}, we conclude that it is sufficient to minimize over elements of $I_v$ such that $\\dist_{H^1}(\\psi_0, \\mathcal{M}_{\\mathcal{E}})=O(v)=\\dist_{L^2}(\\varphi_0,\\mathcal{M}_{\\mathcal{F}})$ for small $v$. Since the $L^2$-projection of $\\varphi_0$ is $\\varphi_{\\rm P}$ by \\eqref{eq:centeringII}, it immediately follows that $\\|\\delta_2\\|_2=O(v)$. We now proceed to show that necessarily also $\\|\\delta_1\\|_{H^1}=O(v)$.\nLet $y',y\\in \\mathbb{R}^{3}$ and $\\theta \\in[0,2\\pi)$ be such that \n\\begin{align}\n\\label{eq:defTheta}\nP_{L^2}^{\\mathcal{M}_{\\mathcal{E}}}(\\psi_0)=\\psi_{\\rm P}^{y'}, \\quad P_{H^1}^{\\mathcal{M}_{\\mathcal{E}}}(\\psi_0)=e^{i\\theta}\\psi_{\\rm P}^y,\n\\end{align}\nwhere we recall that the $L^2$-projection is assumed to be positive by \\eqref{eq:centeringI}. Combining the upper bound derived in the first step with \\cite[Eq.~(53)]{feliciangeli2021persistence}, we get\n\\begin{align}\n\\| \\varphi_0 - \\varphi_{\\rm P}^y \\|_2^2 \\leq C \\left( \\mathcal{G} (\\psi_0, \\varphi_0) - e_{\\rm P} \\right)\\leq Cv^2.\n\\end{align}\nThere exist $\\delta, C_1, C_2>0$ such that \n\\begin{align}\n\\| \\varphi_{\\rm P} - \\varphi_{\\rm P}^y \\|_2\\geq\\begin{cases}\nC_1\\vert y \\vert \\| \\nabla \\varphi_{\\rm P} \\|_2, & |y|\\leq \\delta\\\\\nC_2 & |y|>\\delta\n\\end{cases},\n\\end{align}\nand this allows to conclude that $|y| = O(v)$. In other words, assuming centering w.r.t. to translations in the phonon coordinate (i.e. \\eqref{eq:centeringII}) forces, at low energies, also the centering w.r.t. translations in the electron coordinate, at least approximately. At this point, it is also easy to verify that $\\theta=O(v)$ (and, as an aside, that $|y'|=O(v)$), since, by the upper bound derived in the first step and the coercivity of $\\mathcal{E}$, we have\n\\begin{align}\n\\|\\psi_{\\rm P}^{y'}-e^{i\\theta} \\psi_{\\rm P}^y\\|_2\\leq \\|\\psi_{\\rm P}^{y'}-\\psi_0\\|_{2}+\\|e^{i\\theta} \\psi_{\\rm P}^y-\\psi_0\\|_{2}=O(v).\n\\end{align}\nIn particular, we conclude that\n\\begin{align}\\label{aa}\n\\| \\delta_1 \\|_{H^1} \\leq \\| e^{i\\theta}\\psi_{\\rm P}^{y} - \\psi_0 \\|_{H^1}+ \\| \\psi_{\\rm P} - e^{i\\theta}\\psi_{\\rm P}^{y} \\|_{H^1} = O(v).\n\\end{align} \n\nUsing again \\eqref{eq:expressionG} and that $Q(H_{\\rm P}-4X_{\\rm P})Q\\geq 0$, we conclude that for any $(\\psi_0,\\varphi_0)\\in I_v$ satisfying \n\\eqref{eq:centeringI} and \\eqref{eq:centeringII}, as well as $\\mathcal{G}(\\psi_0,\\varphi_0) \\leq e_{\\rm P} + O(v^2)$, we have\n\\begin{align}\n\\label{eq:firstLB}\n\\mathcal{G}(\\psi_0,\\varphi_0)\\geq e_{\\rm P}+\\expval{H_{\\rm P}}{\\Im\\psi_0}+\\|\\Im\\varphi_0\\|_2^2+O(v^3).\n\\end{align}\nBy arguing as in \\eqref{eq:defVph2}, we see that the conditions $\\text{X}_{\\text{ph}}(0)=0$ and $\\text{V}_{\\text{ph}}(0)=v$ imply that \n\\begin{align}\nP_{\\nabla \\varphi_{\\rm P}} (\\Im \\varphi_0 + v\\alpha^2 \\partial_1 \\Re\\varphi_0)=0 ,\n\\end{align}\nwhere $P_{\\nabla \\varphi_{\\rm P}}$ denotes the projection onto the span of $\\partial_j \\varphi_{\\rm P}$, $1\\leq j\\leq 3$. \nSince $P_{\\nabla \\varphi_{\\rm P}} \\partial_1$ is a bounded operator, and $\\|\\delta_2\\|_2=O(v)$, we find \n\\begin{equation}\\label{325}\n\\|\\Im\\varphi_0\\|_2^2 \\geq \\| P_{\\nabla \\varphi_{\\rm P}} \\Im\\varphi_0\\|_2^2 = v^2 \\alpha^4 \\| \\partial_1 \\varphi_{\\rm P} + P_{\\nabla \\varphi_{\\rm P}} \\partial_1 \\Re \\delta_2 \\|_2^2 \\geq v^2 \\alpha^4 \\| \\partial_1 \\varphi_{\\rm P}\\|_2^2 - O(v^3) .\n\\end{equation}\nWe are left with giving a lower bound on $\\expval{H_{\\rm P}}{\\Im\\psi_0}$, \nunder the condition that \n\\begin{align}\n2\\expval{-i\\nabla}{\\psi_0}=4\\bra{\\Im \\psi_0}\\ket{\\nabla\\Re\\psi_0}=v(1,0,0).\n\\end{align}\nWe already argued in \\eqref{aa} that $\\|\\psi_0-\\psi_{\\rm P}\\|_{H^1}=O(v)$, and therefore\n\\begin{align}\n\\label{eq:constranit-Im}\n4\\bra{\\Im \\psi_0}\\ket{\\nabla\\psi_{\\rm P}}=v(1,0,0)+O(v^2).\n\\end{align}\nCompleting the square, we find \n\\begin{align}\n\\label{eq:completing-square}\n\\expval{H_{\\rm P}}{\\Im\\psi_0} & = \\expval{H_{\\rm P}^{-1}}{H_{\\rm P}\\Im\\psi_0 - v \\partial_1 \\psi_{\\rm P}} \\notag \\\\\n& \\quad + 2v \\bra{\\Im\\psi_0}\\ket{\\partial_1 \\psi_{\\rm P}} - v^2 \\expval{H_{\\rm P}^{-1}}{ \\partial_1 \\psi_{\\rm P}} \\notag \\\\\n& \\geq 2v \\bra{\\Im\\psi_0}\\ket{\\partial_1 \\psi_{\\rm P}} - v^2 \\expval{H_{\\rm P}^{-1}}{ \\partial_1 \\psi_{\\rm P}} .\n\\end{align}\nFrom the constraint \\eqref{eq:constranit-Im} and \\eqref{eq:sandwich}, it thus follows that\n\\begin{align}\\label{329}\n\\expval{H_{\\rm P}}{\\Im\\psi_0} &\\geq \\frac{v^2}{4} + O(v^3).\n\\end{align}\nBy combining \\eqref{eq:firstLB}, \\eqref{325} and \\eqref{329}, we arrive at the final lower bound \n\\begin{align}\nE(v) \\geq e_{\\rm P}+ v^2\\left(\\frac 1 4 + \\alpha^4 \\|\\partial_1\\varphi_{\\rm P}\\|_2^2\\right)+O(v^3) .\n\\end{align}\nAgain, since $\\varphi_{\\rm P}$ is radial, this is of the desired form, and hence the proof is complete.\n\\hfill\\qed\n\n\n\\section{Further Considerations}\n\\label{sec:approaches}\n\nIn this section, we carry out the details related to Remarks \\ref{rem:TW} and \\ref{rem:AlternatePos}. \n\n\\subsection{Effective mass through traveling wave solutions}\n\\label{sec:TW}\n\nA possible way of formalizing the derivation of the effective mass in \\cite{alexandrov2010advances,landau1948effective} relies on traveling wave solutions of the Landau--Pekar equations. \nA traveling wave of velocity $v\\in \\mathbb{R}^{3}$\nis a solution $(\\psi_t, \\varphi_t)$ of \\eqref{eq:LP}\nof the form \n\t\\begin{align}\n\t\\label{eq:TWdef}\n\t(\\psi_t, \\varphi_t) = (e^{-i e_v t} \\psi_v^{\\text{TW}} ( \\cdot - vt ), \\varphi_v^{\\text{TW}} ( \\cdot - vt ) ) \n \\end{align} \n for all $t\\in \\mathbb{R}$, with $e_v\\in \\mathbb{R}$ defining a suitable phase factor. As before, by rotation invariance we can restrict our attention to velocities of the form $v (1,0,0)$ with $v\\in \\mathbb{R}$ in the following.\n \nNote that in the case $\\alpha =0$, where $\\varphi_t = - \\sigma_{\\psi_t}$ for all $t\\in\\mathbb{R}$, the LP equations simplify to a non-linear Schr\\\"odinger equation (also known as Choquard equation). In this case, a traveling wave is given by $\\psi_v^{\\text{TW}} (x)= e^{i x_1 v \/2 } \\psi_{\\rm P} ( x)$ with $e_v= \\mu_{\\rm P} + \\frac{v^2}{4} $, and its energy can be computed to be \n\\begin{align}\n\\mathcal{G}\\left( \\psi_v^{\\text{TW}}, -\\sigma_{\\psi_v^{\\text{TW}}} \\right) = e_{\\rm P} + \\frac{v^2}{4}, \\label{eq:Ev-expansion-alpha-0}\n\\end{align}\nyielding an effective mass $m=1\/2$ at $\\alpha=0$. For the case $\\alpha > 0$, on the other hand, we conjecture that there are no traveling wave solutions of the form \\eqref{eq:TWdef}. \n \n \\begin{Conj}\n For $\\alpha>0$, there are no solutions to the LP equations \\eqref{eq:LP} of the form \\eqref{eq:TWdef} with $v\\neq 0$.\n \\end{Conj}\n \n If one {\\em assumes} the existence of traveling wave solutions, at least for small $v$, one can predict an effective mass that agrees with our formula \\eqref{eq:OurEffMass}, as we shall now demonstrate.\n\nFrom the LP equations \\eqref{eq:LP} one easily sees that a traveling wave solution needs to satisfy \n\\begin{equation}\\label{eq:LPTW}\n\\begin{aligned}\n-i v \\partial_1 \\psi_v^{\\text{TW}} & =\\left( h_{ \\varphi_v^{\\text{TW}}} + e_v \\right) \\psi_v^{\\text{TW}} \\\\\n-i \\alpha^2 v \\partial_1 \\varphi_v^{\\text{TW}} & = \\varphi_v^{\\text{TW}} + \\sigma_{\\psi_v^{\\text{TW}}} .\n\\end{aligned}\n\\end{equation}\nWe shall denote by $E^{\\text{TW}}(v)$ the energy of the traveling wave as a function of the velocity $v\\in \\mathbb{R}$, i.e. \n\\begin{align}\n\\label{eq:Evdef}\nE^{\\text{TW}}(v):=\\mathcal{G}(\\psi_v^{\\text{TW}},\\varphi_v^{\\text{TW}}) .\n\\end{align}\nIn the following, we assume that $e_v=\\mu_{\\rm P}+O(v^2)$ and that the traveling wave is of the form \n\\begin{align}\n\\label{eq:ansatz_TW}\n(\\psi_v^{\\text{TW}}, \\varphi_v^{\\text{TW}}) = \\left( \\frac{\\psi_{\\rm P} + v \\xi_v}{\\| \\psi_{\\rm P} + v \\xi_v \\|_2}, \\varphi_{\\rm P} + v \\eta_v \\right) ,\n\\end{align} \nwith both $\\xi_v$ and $\\eta_v$ bounded in $v$ and converging to some $(\\xi,\\eta)$ as $v\\to 0$. In other words, we assume that the traveling waves have a suitable differentiability in $v$, at least for small $v$, and converge to the standing wave solution $(e^{-i \\mu_{\\rm P} t}\\psi_{\\rm P}, \\varphi_{\\rm P})$ for $v=0$. W.l.o.g. we may also assume that $\\xi_v$ is orthogonal to $\\psi_{\\rm P}$. \n\nWe can then use that \n\\begin{align}\n\\label{eq:exp_norm}\n\\frac{1}{\\| \\psi_{\\rm P} + v \\xi_v \\|_2^2} &= 1- v^2 \\frac{ \\| \\xi_v \\|_2^2}{\\| \\psi_{\\rm P} + v \\xi_v\\|_2^2}\n= 1 - v^2 \\| \\xi \\|_2^2 + o (v^2) \n\\end{align}\nin order to linearize the traveling wave equations \\eqref{eq:TWdef}, obtaining that $(\\xi, \\eta)$ solves\n\\begin{align}\n\\label{eq:xieta}\n\\begin{pmatrix}\ni \\partial_1 \\psi_{\\rm P} \\\\ i \\alpha^2 \\partial_1 \\varphi_{\\rm P} \n\\end{pmatrix}= \n\\begin{pmatrix}\nH_{\\rm P} & 2(2\\pi)^{3\/2}\\psi_{\\rm P}(-\\Delta)^{-1\/2} \\Re\\\\ 2(2\\pi)^{3\/2}(-\\Delta)^{-1\/2} \\psi_{\\rm P} \\Re & 1\n\\end{pmatrix}\n\\begin{pmatrix}\n\\xi \\\\ \\eta\n\\end{pmatrix},\n\\end{align}\nwhere $H_{\\rm P}=h_{\\varphi_{\\rm P}}-\\mu_{\\rm P}$, as defined in \\eqref{eq:def-HP}. Splitting into real and imaginary parts, we equivalently find \n\\begin{align}\n&H_{\\rm P} \\Im \\xi = \\partial_1 \\psi_{\\rm P} \\label{eq:ne_Imxi}\\\\\n&\\Im \\eta =\\alpha^2 \\partial_1 \\varphi_{\\rm P} \\label{eq:ne_Imeta} \\\\\n&H_{\\rm P} \\Re \\xi + 2 (2 \\pi)^{3\/2} \\psi_{\\rm P} \\left( - \\Delta \\right)^{-1\/2} \\Re \\eta = 0 \\label{eq:ne_Rexi}\\\\\n&2 (2 \\pi)^{3\/2} \\left( -\\Delta \\right)^{-1\/2} \\psi_{\\rm P} \\Re \\xi + \\Re \\eta =0. \\label{eq:ne_Reeta}\n\\end{align}\nCombining \\eqref{eq:ne_Rexi} and \\eqref{eq:ne_Reeta} gives $(H_{\\rm P} - 4 X_{\\rm P}) \\Re \\xi =0$, with $X_{\\rm P}$ defined after \\eqref{36}. It was shown in \\cite{lenzmann2009uniqueness} that the kernel of $H_{\\rm P} - 4 X_{\\rm P}$ is spanned by $\\nabla \\psi_{\\rm P}$, hence $\\Re \\xi\\in \\spn \\{\\nabla \\psi_{\\rm P}\\}$. Eq, \n\\eqref{eq:ne_Reeta} then implies that $\\Re \\eta\\in \\spn \\{ \\nabla \\varphi_{\\rm P}\\}$. \n\nUsing these equations and \\eqref{eq:exp_norm} in the expansion \\eqref{eq:expressionG}, it is straightforward to obtain\n\\begin{align}\n\\label{eq:EvexpansionTW}\nE^{\\text{TW}}(v) = e_{\\rm P} + v^2 \\left( \\frac 14 + \\alpha^4 \\| \\partial_1 \\varphi_{\\rm P} \\|_2^2 \\right) + o( v^2), \n\\end{align}\nwhich agrees with \\eqref{eq:Ev-expansion-alpha-0} for the case $\\alpha =0$, and also with \\eqref{eq:thm} to leading order in $v$. In particular, \\eqref{eq:mass-predTW} holds.\n\n\\subsection{Effective mass with alternative definition for the electron's velocity} \n\\label{sec:diff-vel}\nIn this Section, we discuss a different approach to the definition of the effective mass. This approach is based on an alternative way of defining the electron's position and velocity. While in Section \\ref{sec:vel} we use the standard definition from quantum mechanics, here we use a definition similar to the one of the phonons' position and velocity \\eqref{eq:defXph}. For this purpose, we recall Remark \\ref{rem:E} and that $\\delta^*$ has been chosen such that the condition $\\mathcal{E} (\\psi_0) \\leq \\mathcal{G}( \\psi_0, \\varphi_0) \\leq e_{\\rm P} + \\delta^*$ ensures that for all times there exists a unique $L^2$-projection $e^{i\\theta(t)}\\psi_{\\rm P}^{y(t)}$ of $\\psi_t$ onto the manifold $\\mathcal{M}_\\mathcal{E}$. Then, we define the electron's position and velocity by\n\\begin{align}\n\\label{eq:def-tildeX}\n\\widetilde{\\text{X}}_{\\text{el}}(t) = y(t), \\quad \\widetilde{\\text{V}}_{\\text{el}}(t) = \\dot{y}(t).\n\\end{align}\nSimilarly to the conditions (i) and (ii) in Section \\ref{sec:vel}, we define the set of admissible initial data as \n\\begin{align}\n\\widetilde{I}_v = \\lbrace (\\psi_0, \\varphi_0) \\, \\vert \\, \\text{(i),(ii') are satisfied} \\rbrace\n\\end{align}\nwhere \n\\begin{itemize}\n\t\\item[(ii')] $\\widetilde{\\text{V}}_{\\text{el}}(t) =\\text{V}_{\\text{ph}}(0)=v (1,0,0)$. \n\\end{itemize}\n\nNote that we are leaving the parameter $\\dot{\\theta}(0)$ free, which in this case is also relevant. In other words, we have \n\\begin{align}\n\\widetilde{I}_v = \\cup_{\\kappa\\in \\mathbb{R}} \\widetilde{I}_{v,\\kappa},\n\\end{align} \nwhere\n\\begin{align}\n\\widetilde{I}_{v,\\kappa}=\\lbrace (\\psi_0, \\varphi_0) \\, \\vert \\, \\text{(i),(ii') are satisfied and}\\; \\dot{\\theta}(0)=\\kappa \\rbrace.\n\\end{align}\nMinimizing now the energy over all states of the set $\\widetilde{I}_v$\n\\begin{align}\n\\widetilde{E}(v):=\\inf_{(\\psi_0,\\varphi_0)\\in \\widetilde{I}_v} \\mathcal{G}(\\psi_0,\\varphi_0),\n\\end{align}\nleads to an energy expansion in $v$ that differs from the one of Theorem \\ref{thm:freephase} in its second order. \n\n\\begin{prop}\n\t\\label{rmk:EV}\n\tAs $v \\rightarrow 0$, we have \n\t\\begin{align}\n\t\\label{eq:Evtilde}\n\t\\widetilde{E}(v) =e_{\\rm P}+v^2 \\left( \\frac{\\| \\nabla \\psi_{\\rm P} \\|_2^4}{3 \\|\\nabla \\varphi_{\\rm P} \\|_2^2} + \\frac{\\alpha^4}{3} \\| \\nabla \\varphi_{\\rm P} \\|_2^2 \\right) + O(v^3).\t\n\t\\end{align}\n\\end{prop}\n\nThe energy expansion in \\eqref{eq:Evtilde} leads to the effective mass \n\\begin{align}\n \\widetilde{m}_{\\text{eff}} = \\lim_{v \\rightarrow 0} \\frac{\\widetilde{E}(v) -e_{\\rm P}}{v^2\/2} = \\frac{2 \\| \\nabla \\psi_{\\rm P} \\|_2^4}{3 \\|\\nabla \\varphi_{\\rm P} \\|_2^2 } + \\frac{2\\alpha^4}{3} \\| \\nabla \\varphi_{\\rm P} \\|_2^2 \n\\end{align} \nwhich agrees with \\eqref{eq:OurEffMass} and \\eqref{eq:mass-predTW} in leading order for large $\\alpha$ only (and thus still confirms the Landau--Pekar prediction \\cite{landau1948effective}), but differs in the $O(1)$ term. In fact, it turns out that $\\widetilde{m}_{\\text{eff}} < m_{\\text{eff}}$ with $m_{\\text{eff}}$ defined in \\eqref{eq:mass-predTW}. \n\nThis follows from the observation that the trial state\n\\begin{align}\n\\label{eq:trialstate-mass}\n(\\widetilde{\\psi}_0,\\widetilde{\\varphi}_0)=\\left(\\frac{f_v \\psi_{\\rm P}+ivH_{\\rm P}^{-1}\\partial_1 \\psi_{\\rm P}}{\\|f_v \\psi_{\\rm P}+ivH_{\\rm P}^{-1}\\partial_1 \\psi_{\\rm P}\\|}, \\varphi_{\\rm P}+iv\\alpha^2 \\partial_1 \\varphi_{\\rm P}\\right)\n\\end{align}\nwith $f_v = \\frac{1}{2}\\left( 1 + \\sqrt{1 - v^2\/(4 \\| \\partial_1 \\psi_{\\rm P} \\|_2^2)} \\right)$ (which coincides up to terms of order $v^2$ with the trial state \\eqref{eq:trialstate}) is an element of $\\widetilde{I}_{v,\\bar \\kappa}$ for $\\bar \\kappa= - \\mu_{\\rm P}+4\\|\\partial_1 \\psi_{\\rm P}\\|_2^2(f_v-1)$ and is such that $\\mathcal{G} (\\widetilde{\\psi}_0,\\widetilde{\\varphi}_0) = e_{\\rm P} + m_{\\text{eff}}\\; v^2\/2 + O(v^3)$. Thus, $\\widetilde{m}_{\\text{eff}} \\leq m_{\\text{eff}}$ and equality holds if and only if equality (up to terms $o(v^2)$) holds in \\eqref{eq:completing-square2}. This is the case if and only if \n\\begin{align} \nQ_{\\psi_{\\rm P}} \\left( \\Im\\widetilde{\\psi}_0 - cv \\partial_1 \\psi_{\\rm P}\\right) = o(v).\n\\end{align}\nUsing \\eqref{eq:trialstate-mass}, equality holds if and only if\n\\begin{align} \n 0 = H_{\\rm P}^{-1}\\partial_1 \\psi_P - c \\partial_1 \\psi_{\\rm P} = -\\left( x_1\/2 + c \\partial_1 \\right) \\psi_{\\rm P} \n\\end{align}\ni.e. , recalling the radiality of $\\psi_{\\rm P}$, if and only if $\\psi_{\\rm P}$ is a Gaussian with variance $\\sigma^2 = 1\/(2c)$. Since $\\psi_{\\rm P}$ satisfies the Euler--Lagrange equation\n\\begin{align}\nH_{\\rm P} \\psi_{\\rm P}=0 \\iff V_{\\varphi_{\\rm P}} \\psi_{\\rm P}=(-\\Delta+\\mu_P) \\psi_{\\rm P},\n\\end{align}\nit cannot be a Gaussian and therefore $\\tilde{m}_{\\text{eff}}0,\\\\\n u(\\mathbf{x},t=0)=u_0(\\mathbf{x}),&\\mathbf{x}\\in\\mathcal{D},\n\\end{cases}\n\\end{equation}\nwhere $\\epsilon\\geq 0$. We propose an Eulerian-Lagrangian Runge-Kutta finite volume (EL-RK-FV) scheme utilizing weighted essentially non-oscillatory (WENO) schemes with adaptive order (WENO-AO) for spatial reconstruction. The proposed method is designed for one-dimensional problems of the form \\eqref{CDeqn} and extended to higher-dimensional problems via dimensional splitting.\\\\\n\\indent Eulerian-Lagrangian (EL) and semi-Lagrangian (SL) schemes \\cite{Celia1990,Russell2002,Xiu2001} have proven to be computationally effective when solving hyperbolic problems because of their ability to employ high spatial resolution schemes while admitting very large CFL with numerical stability. Generally speaking, an EL or SL method involves working on a background grid so that high-order spatial resolutions can be used (the Eulerian part), and tracing characteristics of the cell boundaries backward\/forward in time to relax the CFL constraint (the Lagrangian part). Such methods have been developed in a wide variety of frameworks: discontinuous Galerkin \\cite{Cai2017,Ding2020,Qiu2011b,Rossmanith2011}, finite difference \\cite{Carrillo2007,Chen2021,Huot2003,Li2022,Qiu2010,Qiu2011,Xiong2019}, finite volume \\cite{Abreu2017,Benkhaldoun2015,Crouseilles2010,Filbet2001,Huang2012,Huang2017}. Although the SL and EL frameworks are similar in spirit, the SL framework assumes exact characteristic tracing and hence poses difficulties when considering nonlinear problems. Two other methods similar to EL and SL methods are the arbitrary Lagrangian-Eulerian (ALE) methods where an arbitrary mesh velocity not necessarily aligned with the fluid velocity is defined \\cite{Boscheri2014,Boscheri2013,Boscheri2015,Boscheri2017,Donea2004,Hirt1974,Peery2000}, and moving mesh methods where the PDE is first evolved in time and then followed by some mesh-redistribution procedure \\cite{Li2001,Li1997,Luo2019,Russell2002,Stockie2001,Tang2003}. The main difference between these two methods and the previously mentioned methods is that they move the mesh adaptively to focus resolving the solution around sharp transitions; whereas EL and SL methods evolve the equation by following characteristics.\\\\\n\\indent Our goal in this paper is to develop a new high-order EL method in the finite volume framework using method-of-lines (MOL) RK time discretizations. Finite volume methods are attractive since they are naturally mass conservative, easy to physically interpret, and modifiable for nonuniform grids. Similar to the recent developments made by Huang, Arbogast, and Qiu \\cite{Huang2012,Huang2017}, we use approximate characteristics to define a traceback space-time region, and then use WENO reconstructions to evaluate the modified flux. Huang and Arbogast developed a re-averaging technique that allows high-order reconstruction of the solution at arbitrary points by applying a standard WENO scheme \\cite{Cockburn1997,Shu2009} over a uniform reconstruction grid that is defined separately. They then used a natural continuous extension \\cite{Zennaro1986} of Runge-Kutta schemes, which requires the solution at several Gaussian nodes of the interval $[t^n,t^{n+1}]$, to evolve the solution along the approximate characteristics. \\\\\n\\indent The novelty of our proposed method is twofold: (1) the partition of space-time regions formed by linear approximations of the characteristic curves, and (2) integrating the differential equation over the partitioned space-time regions, followed by rewriting the space-time integral form of the equation into a spatial-integral time-differential form. In this way, a MOL RK type method can be directly applied for time discretization, thus avoiding the need to use a natural continuous extension of RK schemes. To be more precise, we construct linear approximate characteristics by using the Rankine-Hugoniot jump condition to define the traceback space-time regions. If the linear approximate characteristics are defined with zero velocity, then the proposed EL-RK-FV scheme reduces to the standard RK-FV method. Whereas, when linear space-time curves adequately approximate the exact characteristics, a large time stepping size is still permitted. We use WENO-AO to perform a solution remapping of the uniform cell averages onto the possibly nonuniform traceback cells. The recently developed WENO-AO schemes \\cite{Arbogast2020,Balsara2016,Balsara2020} are robust and guarantee the existence of the linear weights at arbitrary points. We note that Chen, et al. used WENO-AO schemes in the SL framework \\cite{Chen2021}, and Huang and Arbogast have recently used WENO-AO schemes in the Eulerian framework \\cite{Arbogast2019,Arbogast2020}. RK methods are used to evolve the MOL system along the approximate characteristics. Explicit RK methods, such as the strong stability-preserving (SSP) RK methods \\cite{Gottlieb2001}, are used for convection equations; and implicit-explicit (IMEX) RK methods \\cite{Ascher1997,Conde2017,Higueras2014} are used for convection-diffusion equations. In the latter case, the non-stiff convective term is treated explicitly and the stiff diffusive term is treated implicitly. Dimensional splitting is used to extend the one-dimensional algorithm to solve multi-dimensional problems. The proposed method is high-order accurate, capable of resolving discontinuities without oscillations, mass conservative, and stable with large time stepping sizes.\\\\\n\\indent The paper is organized as follows. We discuss the EL-RK-FV algorithm for pure convection problems in Section 2. In Section 3, we discuss the EL-RK-FV algorithm, coupled with IMEX RK schemes, for convection-diffusion equations. Numerical performance of the EL-RK-FV algorithm is shown in Section 4 by applying the algorithm to several linear and nonlinear test problems. Conclusions are made in Section 5. Appendices are listed after the references.\n\n\n\n\n\n\\section{The EL-RK-FV method for pure convection problems}\n\nThe spirit of the EL-RK-FV method is best demonstrated by starting with a pure convection problem in one dimension, i.e., equation \\eqref{CDeqn} with $\\epsilon=0$ and $g(x)=0$. Let the flux be denoted $f(x)$. We first discuss the formulation of the scheme in Section~\\ref{sec2.1}, followed by discussion of high-order spatial reconstruction in Section~\\ref{sec2.2} and time discretization in Section~\\ref{sec2.3}.\n\n\\subsection{\nScheme formulation\n}\n\\label{sec2.1}\n\nWe discretize the spatial domain $[a,b]$ into $N_x$ intervals with $N_x+1$ uniformly distributed nodes\n\\[a=x_{\\frac{1}{2}} 2$. We therefore conclude, that the DIFP can only be stable for $N_2 = N_3 = 1$.\n\nThere are two further exponents that we need to consider to establish the stability of the DIFP solution. Both follow from a scaling relation exploiting the decoupling of the $1$-sector: At the FP, the operators $\\mathcal{O}_1=\\phi_1^2 (\\phi_2^2+\\phi_3^2)$ and $\\mathcal{O}_2 = \\phi_1^2 \\left(\\phi_2^2 - \\frac{N_2}{N_2+N_3}(\\phi_2^2+\\phi_3^2) \\right)$ correspond to eigendirections of the stability matrix. As we know from the Wilson-Fisher FP that the scaling dimension $[\\phi_I^2] = -\\frac{1}{\\nu_I}+d$, we deduce that\n\\begin{equation}\n[\\mathcal{O}_1] = -\\frac{1}{\\nu_1} - \\frac{1}{\\nu_{2+3}} + 2d, \n\\end{equation}\nand accordingly the corresponding critical exponent is given by\n\\begin{equation}\n\\theta_{6} = \\frac{1}{\\nu_{1}} + \\frac{1}{\\nu_{2+3}} - d ~.\n\\end{equation}\nSimilarly, we deduce for the second operator that the corresponding critical exponent is given by\n\\begin{equation}\n\\theta_{4} = \\frac{1}{\\nu_{1}} + y_{2,2} - d ~,\n\\end{equation}\nwhere $y_{2,2}$ is the scaling dimension of the coupling belonging to $\\phi_2^2 - \\frac{N_2}{N_2+N_3}(\\phi_2^2+\\phi_3^2)$ in the two-field case, cf.\\ Ref.\\ \\cite{Eichhorn:2013zza}. The first relation is the one that arises for a two-field DFP, and gives a negative critical exponent (for values of $N_{1} > 1$).\n\nAt fixed $N_{2} = N_{3} = 1$, it is the critical exponent $\\theta_{4}$ that decides about the stability of the FP. Using results from the two-field case (LPA 12 including an anomalous dimension, cf.\\ Ref.\\ \\cite{Eichhorn:2013zza}) to obtain $\\theta_{1} = y_{2,2}$, $\\theta_{2} = \\frac{1}{\\nu_{1}}$, and $\\theta_{3} = \\frac{1}{\\nu_{2+3}}$, we arrive at the results shown in Tab.~\\ref{Tab:DIFP}. The DIFP is the stable FP for $N_{1} \\geq 6$ and $N_{2} = N_{3} = 1$.\n\n\\begin{table}[!t]\n\\renewcommand{\\arraystretch}{1.4}\n\\renewcommand{\\tabcolsep}{4pt}\n\\begin{tabular}{c c c | c c c c c c}\n$N_1$ & $N_2$ & $N_3$& $\\theta_{1}$ & $\\theta_{2}$ & $\\theta_{3}$ & $\\theta_{4}$ & $\\theta_{5}$ & $\\theta_{6}$ \\\\\n\\hline\\hline\n1 & 1 & 1 & 1.765 & 1.571 & 1.459 & 0.336 & -0.042 & 0.030 \\\\ \n2 & 1 & 1 & 1.765 & 1.459 & 1.459 & 0.224 & -0.042 & -0.092 \\\\\n3 & 1 & 1 & 1.765 & 1.367 & 1.459 & 0.132 & -0.042 & -0.174 \\\\\n4 & 1 & 1 & 1.765 & 1.296 & 1.459 & 0.061 & -0.042 & -0.245 \\\\\n5 & 1 & 1 & 1.765 & 1.242& 1.459 & 0.007 & -0.042 & -0.299 \\\\\n6 & 1 & 1 & 1.765 & 1.203 & 1.459 & -0.032 & -0.042 & -0.338 \\\\ \n\\hline\\hline\n\\end{tabular}\n\\caption{\\label{Tab:DIFP}Critical exponents at the DIFP, using the LPA 12 including anomalous dimensions, and employing the above scaling relations.}\n\\end{table}\n\nNote that the derivation of the scaling relations is based on the assumption that the operators corresponding to these couplings are eigenoperators of the stability matrix. This property is, to the best of our knowledge, an assumption in $d = 3$ \\cite{Calabrese:2002bm}, and is usually not true within a truncation of the RG flow. Nevertheless, the stability properties are not incompatible with explicit numerical results within the LPA 8, where the transition to stability occurs already at $N_1=5$. { Explicitly, the critical exponents in LPA 8 at $N_1=5$ and $N_2=N_3=1$ read $\\theta_1=1.783, \\theta_2=1.193, \\theta_3=1.399, \\theta_4=-0.024,\\theta_5=-0.027$ and $\\theta_6=-0.395$.}\n\n\\vskip 5pt\n\n\n\\subsection{Decoupled biconical fixed point}\n\\label{SubSec:Decoupled biconical fixed point}\n\n\\begin{table}[!h]\n\\renewcommand{\\arraystretch}{1.4}\n\\renewcommand{\\tabcolsep}{2.7pt}\n\\begin{tabular}{c c c | c c c c c c}\n$N_1$ & $N_2$ &$N_3$& $\\theta_{1}$ & $\\theta_{2}$ & $\\theta_{3}$ & $\\theta_{4}$ & $\\theta_{5}$ & $\\theta_{6}$ \\\\\n\\hline\\hline\n1 & 1.2 & 1 & {\\bf 1.753} & {\\bf 1.381} & \\textit{1.537} & 0.285 & -0.075 & {\\bf -0.005} \\\\\n1 & 1.4 & 1 & {\\bf 1.535} & {\\bf 1.448} & \\textit{1.537} & 0.105 & -0.015 & {\\bf -0.010} \\\\\n1 & 1.5 & 1 & {\\bf 1.537} & {\\bf 1.462} & \\textit{1.537} & 0.068 & 0.001 & {\\bf -0.001} \\\\ \n1 & 1.2 & 2 & {\\bf 1.753} & {\\bf 1.381} & \\textit{1.399} & 0.161 & -0.200 & {\\bf -0.005} \\\\\n1 & 1.4 & 2 & {\\bf 1.535} & {\\bf 1.448} & \\textit{1.399} & -0.020 & -0.140 & {\\bf -0.010} \\\\\n1 & 1.5 & 2 & {\\bf 1.537} & {\\bf 1.462} & \\textit{1.399} & -0.057 & -0.124 & {\\bf -0.001} \\\\\n1.2 & 1 & 5 & {\\bf 1.753} & {\\bf 1.381} & \\textit{1.193} & -0.051 & -0.413 & {\\bf -0.005} \\\\\n1.5 & 1 & 5 & {\\bf 1.537} & {\\bf 1.462} & \\textit{1.193} & -0.270 & -0.338 & {\\bf -0.001} \\\\\n\\hline\\hline\n\\end{tabular}\n\\caption{\\label{Tab:BFP}Critical exponents at the DBFP in the LPA 8. Critical exponents emphasized in \\textbf{bold} font correspond to those of the two-field BFP, while those in \\textit{italic} arise in the sector with $O(N_3)$ group symmetry. Here, we list the DBFP for values of the field components $N_{I}$, where it is characterized by $\\Delta > 0$ (see text).}\n\\end{table}\n\nBeyond the isotropic and decoupled FP solutions, the two-field models feature another scaling solution, which is the biconical FP \\cite{PhysRevLett.33.813, Calabrese:2002bm, Folk:2008mi}. It is stable only in a restricted parameter region of these models, where $\\left. \\Delta^{\\textrm{2-field}} \\right|_{\\textrm{BFP}} = \\lambda_{2,0} \\lambda_{0,2} - \\lambda_{1,1}^2 > 0$, as it transfers the stability from the IFP to the DFP. Certainly, this FP should similarly manifest itself in the three-field case. While one of the three sectors decouples, the two remaining sectors should feature nondegenerate couplings, and we expect that in a given range of the parameter space such a decoupled BFP will be stable. \n\n{ To obtain as precise results as possible, we should make use of all methods available to us. In fact, results obtained using an $\\epsilon$-expansion around $d=4$ in the two-field case, allow us to infer the stability of the decoupled biconical fixed point in the three field case in one important instance: }\nFrom Refs.\\ \\cite{Folk:2008mi, Calabrese:2002bm}, we know that the biconical FP is stable for $N_1 = 1, N_2 = 2$ in two-field models (and similarly when the sectors are interchanged). Combined with the pattern in Tab.\\ \\ref{Tab:BFP} for the additional critical exponents in the three-field model, we conjecture that the DBFP is stable for $N_1 = 1, N_2 = 2, N_3\\geq 2$ (up to a permutation of the three sectors). To calculate the corresponding critical exponents directly in the three-field model, we expect that an extended truncation will be necessary, taking into account a field-dependent wavefunction renormalization.\n\n\n\\subsection{Search for further stable fixed points}\n\\label{SubSec:Search for further stable fixed points}\n\nWe summarize our results obtained so far in Fig.~\\ref{Fig:StabilityPlot}. The figure shows the stable FP solution for the corresponding values of field components, $(N_1, N_2, N_3)$. Apparently, no stable FP exists in the range $N_{1} < 6$, $N_{2} = N_{3}=1$ (up to a permutation of the sectors) that can be derived from the known scaling solutions in the one- and two-field models. This motivates an independent analysis of fully coupled FPs in the three-field model, which we describe in the following sections.\n\n\\begin{figure}[!t]\n\\includegraphics[width=0.8\\linewidth]{stability_cube_conjectured_modified.pdf}\n\\caption{\\label{Fig:StabilityPlot}(Color online) We show the stable DIFP (large blue dots), the stable DFP (small orange dots) and points without a stable FP (small gray dots) using the LPA 12$+\\eta$ results. We also include points where the DBFP is conjectured to be stable (middle-sized green dots).}\n\\end{figure}\n\n\n\\subsubsection{Stability trading between stable FPs}\n\\label{SubSec:Stability trading between stable FPs}\n\nGenerally, the $\\beta$ functions are non-polynomial functions of a large number of couplings in the multifield models. In the local potential type approximations the number of parameters and couplings increases from 9 to 34 if the order of the truncation is changed from 4 to 8. Thus, finding FPs in the multifield models becomes a highly nontrivial search for zeros of the $\\beta$ functions in a high-dimensional parameter space. In the following, we consider strategies to identify new FP solutions using a simple example. Consider the following $\\beta$ function which is expanded in terms of the coupling $g$ (assuming that higher than quadratic terms are zero):\n\\begin{equation}\n\\beta_g = g ( c + g ),\n\\end{equation}\nwhere $c$ is a function of the parameters of the model (e.g., dimensionality, number of field components, etc.) and possibly other couplings in a given truncation of the theory. Note that such a form captures the essential properties of typical fixed points, as it allows both for a trivial Gaussian FP and a nontrivial interacting FP, as a function of the parameter $c$. The critical exponent at a fixed point is given by\n\\begin{equation}\n\\theta = - \\left. \\frac{\\partial\\beta_g}{\\partial g} \\right|_{\\textrm{FP}} = \\left. - c - 2 g \\right|_{\\rm FP}. \n\\end{equation}\nAssuming that it is the exponent $\\theta$ that decides about the stability of the FP, we may distinguish the following scenarios: For $c < 0$, the interacting FP at $g = -c$ is infrared stable, whereas the Gaussian FP is unstable. As the parameter $c$ increases towards positive values (as a function of, e.g., $N_I$) the two FPs will approach each other. At $c = 0$, both FPs collide and exchange their stability properties. Moving apart again for $c > 0$, the interacting FP has become the unstable one, whereas the noninteracting FP is stable, cf.\\ Fig.\\ \\ref{Fig:StabilityTrading}. This simple example demonstrates that FPs will typically change their stability when they collide, as was also observed in Ref.\\ \\cite{Eichhorn:2013zza} in the two-field model with $O(N_{1})\\oplus O(N_{2})$ symmetry.\n\n\\begin{figure}[!h]\n\\includegraphics[width=0.8\\linewidth]{stability_trading3.pdf}\n\\caption{\\label{Fig:StabilityTrading}To illustrate the mechanism how stability properties are exchanged between FPs, we examine the $\\beta$-function $\\beta = g(c + g)$ for varying parameter $c$ describing the position of the nontrivial FP. As the FPs pass by one another, the local derivatives $\\theta = -\\beta'(g)$ exchange their sign indicating an interchange of their IR stability properties. }\n\\end{figure}\n\n\nInspired by the above stability-trading mechanism\nwe may devise a strategy to identify new FP solutions. Our search for FPs will concentrate on the vicinity of points in coupling space, where a known FP loses its stability.\n\n\n\\subsubsection{Fully coupled fixed points in the three-field model}\n\\label{SubSec:Fully coupled fixed points in the three-field model}\n\nIt turns out that the three-field model works in a different way from the mechanism described above, which applies in the two-field case: Within the LPA 4, we observe that the IFP, DFP, and DIFP have partially overlapping stability regions. No similar behavior occurs in two-field models, where stability regions of different FPs always touch, but never overlap, due to the above stability-trading mechanism. However, this changes dramatically as we include anomalous dimensions: While the IFP inhabits the same points, the DIFP is now only stable for $N_1 \\geq 4$, $N_2 = N_3=1$ etc.\\ Extensive numerical searches did not reveal a stable FP for $N_1 =3$, $N_2 = N_3=1$, and similarly for the cases where the sectors are interchanged. A similar result holds for higher orders of the LPA: As there are a larger number of independent operators that serve as a basis for the LPA, and thus potentially relevant directions in the three-field model, the stability exchange mechanism may not be captured by the simple model considered above. To elucidate the differences in the three-field model, we will focus on results obtained within the LPA to 8th order. \n\n\\begin{figure}[!t]\n\\includegraphics[width=\\linewidth]{BFPDFPcollision.pdf}\n\\caption{\\label{BFPDFPcollision}We plot $\\lambda_{1,0,1}, \\lambda_{0,1,1}$, and $\\lambda_{1,1,0}$ as a function of $N_I$ for the three DBFPs. At $N_I=1.244$ they each pass through the origin of that coordinate system, where the DFP sits. For $N_I>1.244$, the three BFPs move away from each other towards more negative values of the mixed couplings.}\n\\end{figure}\n\n\\begin{figure}[!h]\n\\includegraphics[width=0.8\\linewidth]{BFPDFPthetas.pdf}\n\\includegraphics[width=0.8\\linewidth]{BFPDFPcoords.pdf}\n\\caption{\\label{BFPDFPthetas}We plot the fifth-largest critical exponent at the DBFPs (blue points of increasing value) and the corresponding critical exponent at the DFP (red points of decreasing value) as a function of $N_I$ (upper panel). Below, we show the coordinates of the couplings $\\lambda_{2,0,0}$ and $\\lambda_{0,2,0}$ at the DBFP (blue diamonds and black squares) and the couplings $\\lambda_{2,0,0}= \\lambda_{0,2,0} = \\lambda_{0,0,2}$ at the DFP (red dots) in the vicinity of the collision point.}\n\\end{figure}\n\nAs a first example, let us consider the point $N_1 = N_2 = N_3 \\approx 1.25$. Here, the DFP is stable, but it gains three additional relevant directions around $N_1 = N_2 = N_3 \\approx 1.244$. As within the two-field model, this point is marked by a collision with a decoupled biconical FP (DBFP). The main difference is that within a three-field model, three generalizations of the BFP exist, cf.\\ Sec.\\ \\ref{SubSec:Decoupled biconical fixed point}. For $N_I =1.25$, $I = 1,2,3$, all of these FPs feature five relevant directions. Toward smaller $N_I$, all three DBFPs approach the DFP, and simultaneously collide with it at $N_I \\approx 1.244$. At this point, the DFP gains one relevant direction from each of the three DBFPs, which subsequently feature only four relevant exponents, cf.\\ Fig.\\ \\ref{BFPDFPcollision} and Fig.~\\ref{BFPDFPthetas}. Thus this FP loses stability in a fixed-point collision. The central difference to the two-field case lies in the fact that the symmetry of the model which forces the DFP to collide with three other FPs simultaneously (for $N_1 = N_2 = N_3$ there is an exchange symmetry $N_1 \\leftrightarrow N_2 \\leftrightarrow N_3$). As each of them starts off with five relevant directions, the collision does not produce a stable FP for $N_I \\leq 1.24$, but instead leaves behind three FPs with four relevant directions each.\n\nAs a second example of new behavior in three-field models, we consider the point where the DBFP collides with the DIFP in the LPA 8. We fix $N_2 =1$ and $N_3=4$: Then the DBFP has five positive critical exponents at $N_1=1$. Going to larger values of $N_1$, it collides with the DIFP at $N_1 \\approx 1.6$. During this collision, the DIFP becomes unstable, and the DBFP gains one negative critical exponent, cf.\\ Fig.\\ \\ref{BFPDIFP}. Similar to the previous scenario, this FP collision is not sufficient to make the DBFP stable. Following the DBFP to even larger values of $N_1$, it undergoes another collision, this time being hit by two other FPs simultaneously. This is a novel feature that is not observed in simple two-field models. Starting in the region $N_1 < 1.2$, these FPs do not seem to exist for real values of the couplings -- at least no sign of them showed up in extensive numerical searches. They can be thought of as being created at the collision. Following this collision, they quickly move away from the collision point for increasing values of $N_{1}$. Each of the newly created FPs features four relevant critical exponents, while the DBFP is stable, cf.\\ Fig.\\ \\ref{BFPcollision}. Both new FPs are anisotropic, and define a new universality class that occurs for the first time if three fields are coupled, cf.\\ Tab.\\ \\ref{Tab:AFP}.\n\n\\begin{figure}[!h]\n\\includegraphics[width=0.9\\linewidth]{BFPcoords.pdf}\\\\\n\\includegraphics[width=0.9\\linewidth]{BFPthetas.pdf}\\\\\n\\includegraphics[width=0.9\\linewidth]{DIFPtheta.pdf}\n\\caption{\\label{BFPDIFP}(Color online) Here, we show the couplings $\\lambda_{2,0,0}$, $\\lambda_{0,2,0}$, and $\\lambda_{1,1,0}$ at the DBFP (blue circles, purple squares, cyan diamonds) and $\\lambda_{2,0,0} = \\lambda_{0,2,0} = \\lambda_{1,1,0}$ (red triangles) at the DIFP (uppermost panel). We show the fourth- and fifth-largest critical exponent (middle panel) at the DBFP as a function of $N_1$ for $N_2=1, N_3=4$. Around $N_1 \\approx 1.16$, the three couplings are clearly degenerate, as expected for a collision with the DIFP. At the same point, the fourth critical exponent crosses zero and becomes negative. Simultaneously, the fifth critical exponent at the DIFP (lower panel) crosses zero and becomes positive.}\n\\end{figure}\n\n\\begin{figure}[!h]\n\\includegraphics[width=\\linewidth]{AFPcollision.pdf}\\\\\n\\includegraphics[width=\\linewidth]{AFPtheta.pdf}\\\\\n\\caption{\\label{BFPcollision}Here, we show the couplings $\\lambda_{1,0,1}$ and $\\lambda_{0,1,1}$ at the two anisotropic FPs. As a function of $N_1$ they converge towards zero, which is where the DBFP is sitting. The collision occurs at $N_1 \\approx 1.2$, which is where the fifth critical exponent of these two FPs (which numerically is nearly the same for both) approaches zero.}\n\\end{figure}\n\n\\begin{table}[!t]\n\\renewcommand{\\arraystretch}{1.4}\n\\renewcommand{\\tabcolsep}{2.7pt}\n\\begin{tabular}{ccc|ccc|ccc|c}\n$N_1$ & $N_2$ & $N_3$ & $\\lambda_{2,0,0}$ & $\\lambda_{0,2,0}$ & $\\lambda_{0,0,2}$ & $\\lambda_{1,1,0}$ & $\\lambda_{1,0,1}$ & $\\lambda_{0,1,1}$ & $\\lambda_{1,1,1}$\\\\ \\hline \\hline\n2 & 2 & 2 & 6.2 & 6.4& 6.4 & 1.8 & 1.8 & -2.2 & -2.6\\\\\n2 & 1 & 2 & 6.5 & 7.5 & 6.7 & 1.2 & 0.5 & -2.1 & -1.4\\\\\n2 & 1 & 3 & 6.2 & 7.0 & 5.8 & 3.2 & 0.7 & -1.8 & -2.2\\\\\n2 & 1 & 4 & 6.0 & 6.5 & 5.1 & 4.6 & 0.6 & -1.3 & -1.6\\\\ \\hline \\hline\n\\end{tabular}\n\\vskip 4pt\n\\renewcommand{\\arraystretch}{1.4}\n\\renewcommand{\\tabcolsep}{6.8pt}\n\\begin{tabular}{ccc|ccc|ccc}\n{} & {} & {} & $\\theta_1$ & $\\theta_2$ & $\\theta_3$ & $\\theta_4$ & $\\theta_5$ & $\\theta_6$ \\\\ \\hline \\hline\n2 & 2 & 2 & 1.70 & 1.33 & 1.29 & 0.20 & -0.17 & -0.24\\\\\n2 & 1 & 2 & 1.62 & 1.40 & 1.35 & 0.08 & -0.05 & -0.23\\\\\n2 & 1 & 3 & 1.71 & 1.34 & 1.28 & 0.15 & -0.10 & -0.30\\\\\n2 & 1 & 4 & 1.77 & 1.32 & 1.23 & 0.16 & -0.07 & -0.35\\\\ \\hline \\hline\n\\end{tabular}\n\\caption{\\label{Tab:AFP}We list selected FP values and critical exponents of the first anisotropic FP in the LPA 8, restricting ourselves to couplings at fourth order in the fields, and additionally giving the value of the coupling associated with the three-field operator $\\sim \\phi_1^2 \\phi_2^2 \\phi_3^2$.}\n\\end{table}\n\nDue to the complicated dynamics of FPs in the parameter space of three-field models, it seems that these models do not necessarily feature a stable FP solution for all values of the model parameters. One might now wonder, why FPs of three-field models show such a disproportionate increase in the number of relevant directions for small values of $N_I$. While the Wilson-Fisher FP has one relevant direction, and two-field models always feature FPs with only two relevant directions, three-field models show FPs with more than three relevant directions. We conjecture that this is related to the fact that the transition from a two-field to a three-field model implies the existence of more than one additional new class of operators: While the transition from one to two fields only adds the mixed interactions $\\sim \\left( \\phi_1^{2} \\right)^{m} \\left( \\phi_2^{2} \\right)^{n}$, the transition to three-field models features three different classes of mixed interactions $\\sim \\left( \\phi_I^{2} \\right)^{m} \\left( \\phi_J^{2} \\right)^{n}$, $I,J = 1,2,3$, in addition to the new three-field couplings $\\sim \\left( \\phi_1^{2}\\right)^{l} \\left( \\phi_2^{2} \\right)^{m} \\left( \\phi_3^{2} \\right)^{n}$. These seem to play a particularly important role for small $N_I$ and imply that the IR scaling properties cannot be accounted for by only three free relevant parameters. Note that this does not imply that the corresponding additional couplings need to become relevant. Since these operators do not necessarily correspond to eigendirections of the RG, they may mix with other operators, and yield corrections to the scaling spectrum.\n\nApplying our results to determine the properties of possible multicritical points in phase diagrams for systems with three competing order parameters, we may conclude that models with small $N_I$ will typically feature a first order, rather than a second order (multicritical) transition. In particular, this applies to the phenomenologically relevant model of three interacting $Z_{2}$-Ising fields, i.e., $Z_{2} \\oplus Z_{2} \\oplus Z_{2}$ symmetry.\n\n\n\\section{Summary and conclusions}\n\\label{Sec:Summary and conclusions}\n\nHere, we present a renormalization group study of IR stable FP solutions in three-field models with $O(N_1) \\oplus O(N_2) \\oplus O(N_3)$ symmetry. Our main results regarding the existence of stable FPs are summarized in Fig.~\\ref{Fig:StabilityPlot}. Models in this class exhibit FPs that generalize the Wilson-Fisher FP, falling into three distinct categories, each characterized by the degree of symmetry-enhancement. We find a decoupled FP, a partially isotropic FP solution, and a fully isotropic FP solution. Their scaling spectrum can be deduced partially by considering perturbations around the single- and two-field models with $O(N)$ and $O(N_{1})\\oplus O(N_{2})$ symmetry, respectively. We proceed by deriving scaling relations between different critical exponents to discover the stability properties of nontrivial FPs in the three-field case. Apart from the generalized Wilson-Fisher scaling solutions, we identify a decoupled biconical FP whose scaling properties are partly inherited from the BFP in the $O(N_{1})\\oplus O(N_{2})$ symmetric model. As a main result of this work, we find that these FPs all show a significantly larger number of relevant critical exponents than in the two-field case, in the region of small $N_1$, $N_2$, and $N_3$. We tentatively connect this result to the existence of a large number of mixed interactions, and further conjecture, that similar results will hold for models with $n>3$ interacting fields.\n\nSummarizing our results, we find no IR stable FP for a small number of fields ($N_{1} < 6$, $N_{2} = N_{3} = 1$, up to permutations of the fields) in $d = 3$ dimensions. This result is certainly unexpected, as there is no evidence for similar behavior in coupled two-field models. While, in principle, we cannot exclude the possibility that stable FPs exist in that region of parameter space, we find no evidence for their existence in extensive numerical searches for FPs of the nonperturbative $\\beta$ functions. The identification of FPs in the three-field field models with $O(N_{1})\\oplus O(N_{2})\\oplus O(N_{3})$ is in general a difficult problem, since the search has to proceed through a high-dimensional coupling space. Nevertheless, the understanding of basic stability transitions between different FPs serves as a guiding principle to single out possible candidates for nontrivial FPs. Quite generally, in coupled-field models, stability seems to be inherited from single mergers or collisions of different FPs. Searches around such stability transition points have not yielded any FP that carries over the stability properties from the IFP (stable at small, noninteger values of $N_{I}<1$) to the decoupled FP. This indicates that the dynamics of FPs in $O(N_{1})\\oplus O(N_{2})\\oplus O(N_{3})$ symmetric three-field, or general $\\bigoplus_{I} O(N_{I})$ symmetric multifield models is very different from that encountered in the simpler $O(N_1)\\oplus O(N_2)$-type models.\n\nFrom these results, we may conclude that models with phenomenological relevance such as, e.g., the $Z_{2} \\oplus Z_{2} \\oplus Z_{2}$ symmetric model will not feature a multicritical point in its phase diagram. Certainly, it is challenging to find three parameters that are accessible experimentally, and may be tuned to the multicritical point. This would be necessary to quantify the scaling behavior close to the corresponding FP, or to show the absence of such a transition. Nevertheless, it is conceivable that in the context of ultracold atomic systems such a control of the system might be achievable \\cite{doi:10.1146\/annurev-conmatphys-070909-104059}. \n\nFinally, let us comment on the general applicability of our results to other systems of interest. The renormalization group flow equations are derived for general $d$ Euclidean dimensions and can be applied to $d = 2$, relevant for critical behavior of low-dimensional condensed-matter systems, and $d = 4$, for multifield models of inflation, as well as possible extensions of the Higgs sector of the standard model. We leave this for future work.\n\n\\begin{acknowledgements}\nWe thank P.~Calabrese for discussions.\nResearch at Perimeter Institute is supported by the Government of Canada through Industry Canada and by the Province of Ontario through the Ministry of Research and Innovation. This work is supported by DOE grant No.\\ DE-FG0201ER41195. M.M.S. is supported by the grant ERC-AdG-290623 and DFG grant FOR 723. A.E. would like to thank the University of Heidelberg and the University of Illinois at Chicago for hospitality during the early stages of this work.\n\\end{acknowledgements}\n\n\\begin{appendix}\n\\label{Appendix}\n\n\\section{Functional renormalization group}\n\\label{Sec:Functional renormalization group}\n\nThe nonperturbative functional renormalization group \\cite{Wetterich:1992yh,Berges:2000ew,Polonyi:2001se,Pawlowski:2005xe,Gies:2006wv,Delamotte:2007pf} defines a functional flow for the scale dependent effective action $\\Gamma_{k}$, which interpolates between the microscopic action $S$ defined at some ultraviolet cutoff scale $k = \\Lambda$ and the full effective action $\\Gamma_{k\\rightarrow 0} = \\Gamma$, when the renormalization group parameter $k$ is removed. It is given by\n\\begin{equation}\n\\partial_t \\Gamma_{k} = \\frac{1}{2} \\Tr \\int\\! \\frac{d^{d}q}{(2\\pi)^{d}} \\left(\\Gamma_k^{(2)} (q) + R_k(q) \\right)^{-1} \\partial_t R_k(q) ~,\n\\label{Eq:WetterichEquation}\n\\end{equation}\nwhere the logarithmic scale derivative is written in terms of the parameter $t = \\ln (k\/\\Lambda)$, and the second functional derivative $\\Gamma_k^{(2)}(p,q) = \\frac{\\delta^2 \\Gamma_k}{\\delta \\chi(-p) \\delta \\chi(q)}$, and $\\Gamma^{(2)}(q) (2\\pi)^{d} \\delta^{(d)}(p-q) \\equiv \\Gamma^{(2)} (p,q)$. Here, $\\chi$ denotes the complete field content of our model and the trace $\\Tr$ denotes a summation over internal degrees of freedom, i.e., both fields and field components. The regulator function $R_{k}$ implements a mass-like cutoff and regulates the infrared divergences. We take $R_{I J} (q) = R_{I}(q) \\delta_{I J}$, $I, J = 1,2,3$, while the momentum dependence is given by $R_{I}(q) = Z_{I} (k^{2} - q^{2}) \\theta(k^{2} - q^{2})$, where the wavefunction renormalization is scale dependent. This choice is referred to as the optimized regulator \\cite{Litim:2000ci, Litim:2001up} which allows us to derive fully analytic expressions for the nonperturbative $\\beta$ functions, and it is thus a convenient choice to identify possible scaling solutions.\n\n\\onecolumngrid\n\\section{Scale dependent effective potential}\n\\label{Sec:Scale dependent effective potential}\n\nWe derive the renormalization group flow equation for the effective potential by plugging our \\textit{ansatz} Eq.\\ \\eqref{Eq:EffectivePotential} into the flow equation \\eqref{Eq:WetterichEquation} and projecting the right-hand-side onto a constant field configuration:\n\\begin{eqnarray}\n\\partial_t u_k &=& -d u_k +\\sum_I (d-2+ \\eta_I) \\rho_I \\partial_{\\rho_{I}} u_{k} + { 2 v_{d} \\sum_{I} \\Big\\{ (N_I-1) l_{0}^{(I)} \\left( \\partial_{\\rho_{I}} u_{k} \\right)} \\nonumber\\\\ && { +\\: l_{R}^{(I)} \\left( \\left\\{ \\partial_{\\rho_{J}} u_{k} + 2 \\rho_{J} \\partial_{\\rho_{J}}^{2} u_{k} \\right\\}, \\left\\{ 4 \\rho_{J} \\rho_{K} \\partial_{\\rho_{J}} \\partial_{\\rho_{K}} u_{k} \\right\\} \\right)} \\Big\\} ~. \n\\label{Eq:FlowEquationPotential}\n\\end{eqnarray}\nHere, $v_{d}^{-1} = 2^{d+1} \\pi^{\\frac{d}{2}} \\Gamma\\left(\\frac{d}{2}\\right)$ arises from the volume integration, and the anomalous dimensions are defined as $\\eta_I = - \\partial_t \\ln Z_I$. In the following the notation follows Sec.\\ \\ref{Sec:Model}. The threshold functions $l_{0}^{(I)}$ and $l_{R}^{(I)}$ (see, e.g., \\cite{Braun:2011pp}) define the diagrammatic contributions to the renormalization group flow of the scale dependent effective potential, where the upper index indicates the corresponding sector. Using an optimized regulator function \\cite{Litim:2000ci,*Litim:2001up} the threshold functions take the following form:\n\\begin{eqnarray}\nl_{0}^{(I)}(w_{I}) &=& \\frac{2}{d} \\left( 1 - \\frac{\\eta_{I}}{d+2}\\right)\\frac{1}{1+w_{I}} ~, \\\\\n\\hspace{-20pt} l_{R}^{(1)}(\\{ w_{I} \\}, \\{ \\delta_{I,J}^{2} \\}) &=& \\frac{2}{d} \\left( 1 - \\frac{\\eta_{1}}{d+2}\\right) \\frac{\\left( 1 + w_{2} \\right) \\left( 1 + w_{3} \\right) - \\delta_{2,3}^{2}}{2 \\delta_{1,2} \\delta_{1,3} \\delta_{2,3} \n- \\delta_{1,2}^{2} \\left( 1 + w_{3} \\right) - \\delta_{1,3}^{2} \\left( 1 + w_{2} \\right) - \\delta_{2,3}^{2} \\left( 1 + w_{1} \\right)\n+ \\prod_{J} \\left( 1 + w_{J} \\right)} ~, \\\\\n\\hspace{-20pt} l_{R}^{(2)}(\\{ w_{I} \\}, \\{ \\delta_{I,J}^{2} \\}) &=& \\frac{2}{d} \\left( 1 - \\frac{\\eta_{2}}{d+2}\\right) \n\\frac{\\left( 1 + w_{1} \\right) \\left( 1 + w_{3} \\right) - \\delta_{1,3}^{2}}{2 \\delta_{1,2} \\delta_{1,3} \\delta_{2,3} \n- \\delta_{1,2}^{2} \\left( 1 + w_{3} \\right) - \\delta_{1,3}^{2} \\left( 1 + w_{2} \\right) - \\delta_{2,3}^{2} \\left( 1 + w_{1} \\right)\n+ \\prod_{J} \\left( 1 + w_{J} \\right)} ~, \\\\\n\\hspace{-20pt} l_{R}^{(3)}(\\{ w_{I} \\}, \\{ \\delta_{I,J}^{2} \\}) &=& \\frac{2}{d} \\left( 1 - \\frac{\\eta_{3}}{d+2}\\right) \n\\frac{\\left( 1 + w_{1} \\right) \\left( 1 + w_{2} \\right) - \\delta_{1,2}^{2}}{2 \\delta_{1,2} \\delta_{1,3} \\delta_{2,3} \n- \\delta_{1,2}^{2} \\left( 1 + w_{3} \\right) - \\delta_{1,3}^{2} \\left( 1 + w_{2} \\right) - \\delta_{2,3}^{2} \\left( 1 + w_{1} \\right)\n+ \\prod_{J} \\left( 1 + w_{J} \\right)} ~.\n\\end{eqnarray}\nNote, that the nature of the interactions is such, that only the radial modes $l_{R}^{(I)}$ are affected by the couplings, e.g., $\\delta_{1,2}^{2} = 4 \\kappa_{1} \\kappa_{2} \\lambda_{1,1,0}^{2}$, and equivalent couplings between the remaining sectors. One may easily check that in the limit of vanishing couplings, $\\delta_{1,3} \\rightarrow 0$ and $\\delta_{2,3} \\rightarrow 0$, the radial contributions $l_{R}^{(I)}$ in the $(1,2)$ sectors reduce to the threshold functions that were already derived in the two-field model \\cite{Eichhorn:2013zza}:\n\\begin{eqnarray}\n\\hspace{-20pt} l_{0}^{(I)}( w_{I} ) &=& \\frac{2}{d} \\left( 1 - \\frac{\\eta_{I}}{d+2}\\right)\\frac{1}{1+w_{I}} ~, \n\\label{Eq:ThresholdFunctionL0} \\\\\n\\hspace{-20pt} l_{R}^{(1)}(\\{ w_{I} \\}, \\{ \\delta_{I,J}^{2} \\}) &=& \\frac{2}{d} \\left( 1 - \\frac{\\eta_{1}}{d+2}\\right) \\frac{1 + w_{2}}{\\left( 1 + w_{1} \\right) \\left( 1 + w_{2} \\right)- \\delta_{1,2}^{2}} ~, \\\\\n\\hspace{-20pt} l_{R}^{(2)}(\\{ w_{I} \\}, \\{ \\delta_{I,J}^{2} \\}) &=& \\frac{2}{d} \\left( 1 - \\frac{\\eta_{2}}{d+2}\\right) \\frac{1 + w_{1}}{\\left( 1 + w_{1} \\right) \\left( 1 + w_{2} \\right)- \\delta_{1,2}^{2}} ~.\n\\label{Eq:ThresholdFunctionLTildeR}\n\\end{eqnarray}\nFrom \\eqref{Eq:FlowEquationPotential} the flow equations for the couplings are derived by the differentiation with respect to the fields and successive projection onto a nonvanishing constant background field configuration $\\rho_{I} = \\kappa_{I}$, defined by the minimum of the effective potential. For some FPs it might be necessary to employ an expansion point where one or several of the $\\kappa_{I}$ are vanishing. In this case $\\beta_{\\kappa_{3}} \\equiv 0$. For a detailed discussion of this issue, we refer to Ref.\\ \\cite{Eichhorn:2013zza}. \n\nWe obtain the $\\beta$ functions for the couplings $\\lambda_{l,m,n}$, $l + m + n \\geq 2$:\n\\begin{equation}\n\\beta_{\\lambda_{l,m,n}} = \\left. \\left(\\frac{\\partial^{l+m+n}}{\\partial \\rho_1^{l} \\partial \\rho_2^{m} \\partial \\rho_3^{n}} \\partial_t u_k + \\beta_{\\kappa_1} \\frac{\\partial^{(l+1)+m+n} u_k}{\\partial \\rho_{1}^{l+1} \\partial \\rho_2^m \\partial \\rho_3^n} + \\beta_{\\kappa_2} \\frac{\\partial^{l+(m+1)+n} u_k}{\\partial \\rho_{1}^{l} \\partial \\rho_2^{m+1} \\partial \\rho_3^n} + \\beta_{\\kappa_3} \\frac{\\partial^{l+m+(n+1)} u_k}{\\partial \\rho_{1}^{l} \\partial \\rho_2^m \\partial \\rho_3^{n+1}} \\right) \\right|_{\\rho_I = \\kappa_I} ~,\n\\end{equation}\nwhere $\\beta_{\\lambda_{l,m,n}} \\equiv \\partial_{t} \\lambda_{l,m,n}$. The $\\beta$ functions for the scale dependent dimensionless field expectation values $\\kappa_{I}$, $I = 1,2,3$, are given by\n\\begin{eqnarray}\n\\beta_{\\kappa_1} &=& \\frac{1}{\\Delta} \\left\\{\n- \\Delta_{2,3} \\partial_{\\rho_1} \\partial_t u_k + \\left(\\lambda_{1,1,0} \\lambda_{0,0,2} - \\lambda_{1,0,1} \\lambda_{0,1,1}\\right) \\partial_{\\rho_2} \\partial_t u_k \\right. \\nonumber \\\\ && \\left.\\left. +\\: \\left(\\lambda_{1,0,1} \\lambda_{0,2,0} - \\lambda_{1,1,0} \\lambda_{0,1,1}\\right) \n\\partial_{\\rho_3} \\partial_t u_k \\right\\} \\right|_{\\rho_I = \\kappa_I} ~, \n\\end{eqnarray}\n\\begin{eqnarray}\n\\beta_{\\kappa_2} &=& \\frac{1}{\\Delta} \\left\\{ - \\Delta_{1,3} \\partial_{\\rho_2} \\partial_t u_k + \\left( \\lambda_{2,0,0} \\lambda_{0,1,1} - \\lambda_{1,1,0} \\lambda_{1,0,1} \\right) \\partial_{\\rho_3} \\partial_t u_k \\right. \\nonumber\\\\ \n&& \\left.\\left. +\\: \\left( \\lambda_{0,0,2} \\lambda_{1,1,0} - \\lambda_{1,0,1} \\lambda_{0,1,1} \\right) \\partial_{\\rho_1} \\partial_t u_k \\right\\} \\right|_{\\rho_I = \\kappa_I} ~, \\\\\n\\beta_{\\kappa_3} &=& \\frac{1}{\\Delta} \\left\\{\n- \\Delta_{1,2} \\partial_{\\rho_3} \\partial_t u_k + \\left( \\lambda_{2,0,0} \\lambda_{0,1,1} - \\lambda_{1,1,0} \\lambda_{1,0,1} \\right) \\partial_{\\rho_2} \\partial_t u_k \\right. \\nonumber\\\\ && \\left.\\left. +\\: \\left( \\lambda_{0,2,0} \\lambda_{1,0,1} - \\lambda_{1,1,0} \\lambda_{0,1,1}\\right) \\partial_{\\rho_1} \\partial_t u_k \\right\\} \\right|_{\\rho_I = \\kappa_I} ~. \n\\end{eqnarray}\nHere, we have defined the coupling parameter\n\\begin{equation}\n\\Delta_{1,2} = \\lambda_{2,0,0} \\lambda_{0,2,0} - \\lambda_{1,1,0}^{2} ~,\n\\end{equation}\nand equivalently $\\Delta_{1,3}$ and $\\Delta_{2,3}$ (defined from the remaining quartic couplings in the three-field model), as well as the parameter\n\\begin{eqnarray}\n\\Delta &=& - 2 \\left( \\lambda_{2,0,0} \\lambda_{0,2,0} \\lambda_{0,0,2} - \\lambda_{1,1,0} \\lambda_{1,0,1} \\lambda_{0,1,1} \\right) + \\lambda_{0,0,2} \\Delta_{1,2} + \\lambda_{0,2,0} \\Delta_{1,3} + \\lambda_{2,0,0} \\Delta_{2,3} ~.\n\\end{eqnarray}\nThese parameters quantify the symmetry enhancement properties of the system. In particular, for certain symmetry enhanced FPs, these quantities vanish exactly.\n\n\\section{Wavefunction renormalization and anomalous dimensions}\n\\label{Sec:Wavefunction renormalization and anomalous dimensions}\n\nTo determine the scale dependence of the field independent renormalization factor $Z_{I}$ from the functional flow equation \\eqref{Eq:WetterichEquation} we perform a projection of the flow onto operators of the type $\\sim \\left( \\partial_{\\mu} \\phi_{I} \\right)^{2} $. This yields the scale dependence of the coefficient for the corresponding operator in the effective action. We have\n\\begin{equation}\n\\partial_{t} Z_{I} = \\frac{(2\\pi)^{d}}{\\delta^{(d)}(0)} \\lim_{Q\\rightarrow 0} \\frac{\\partial}{\\partial Q^{2}} \\frac{\\delta^{2}}{\\delta \\varphi_{I}(-Q) \\delta \\varphi_{I}(Q)} \\partial_{t} \\Gamma_{k} ~,\n\\label{Eq:RenormalizationFactorProjection}\n\\end{equation}\nwhere the functional derivatives are taken with respect to the Nambu-Goldstone (NG) degrees of freedom $\\varphi_{I}$ in the $I$-sector. Note, there is no summation implied over the $I$-index. For details of the derivation in the context of the $O(N)$ vector model, we refer to \\cite{Tetradis:1993ts,Berges:2000ew}. \n\nThe anomalous dimensions are defined via the scaling contribution to the wavefunction renormalization, i.e., $\\eta_{I} = - \\partial_{t} \\ln Z_{I}$, and take the following form in the three-field model:\n\\begin{eqnarray}\n\\eta_{1} &=& \\frac{16 v_{d}}{d} \\Xi^{-1} \\left\\{ \\kappa_{2} \\lambda_{1,1,0}^2 + \\kappa_{3} \\lambda_{1,0,1}^2 + \\kappa_{1} \\left( \\lambda_{2,0,0} + 2 \\kappa_{2} \\Delta_{1,2} + 2 \\kappa_{3} \\Delta_{1,3} + 4 \\kappa_{2} \\kappa_{3} \\Delta \\right)^2 \\right. \\nonumber\\\\ && \\left. +\\: 4 \\kappa_{2} \\kappa_{3} \\left( \\Delta_{2,3} \\left( \\lambda_{2,0,0} + \\kappa_{2} \\Delta_{1,2} + \\kappa_{3} \\Delta_{1,3} \\right) - \\Delta \\left(1 + \\kappa_{2} \\lambda_{0,2,0} + \\kappa_{3} \\lambda_{0,0,2} \\right) \\right) \\right\\} ~, \n\\label{Eq:Eta1} \\\\\n\\eta_{2} &=& \\frac{16 v_{d}}{d} \\Xi^{-1} \\left\\{ \\kappa_{1} \\lambda_{1,1,0}^{2} + \\kappa_{3} \\lambda_{0,1,1}^2 + \\kappa_{2} \\left( \\lambda_{0,2,0} + 2 \\kappa_{1} \\Delta_{1,2} + 2 \\kappa_{3} \\Delta_{2,3} + 4 \\kappa_{1} \\kappa_{3} \\Delta \\right)^2 \\right. \\nonumber\\\\ \n&& \\left. +\\: 4 \\kappa_{1} \\kappa_{3} \\left( \\Delta_{1,3} \\left( \\lambda_{0,2,0} + \\kappa_{1} \\Delta_{1,2} + \\kappa_{3} \\Delta_{2,3} \\right) - \\Delta \\left(1 + \\lambda_{2,0,0} \\kappa_{1} + \\lambda_{0,0,2} \\kappa_{3} \\right) \\right) \\right\\} ~, \n\\label{Eq:Eta2} \\\\\n\\eta_{3} &=& \\frac{16 v_{d}}{d} \\Xi^{-1} \\left\\{ \\kappa_{1} \\lambda_{1,0,1}^2 + \\kappa_{2} \\lambda_{0,1,1}^2 + \\kappa_{3} \\left( \\lambda_{0,0,2} + 2 \\kappa_{1} \\Delta_{1,3} + 2 \\kappa_{2} \\Delta_{2,3} + 4 \\kappa_{1} \\kappa_{2} \\Delta \\right)^2 \\right. \\nonumber\\\\ \n&& \\left. +\\: 4 \\kappa_{1} \\kappa_{2} \\left( \\Delta_{1,2} \\left( \\lambda_{0,0,2} + \\kappa_{1} \\Delta_{1,3} + \\kappa_{2} \\Delta_{2,3} \\right) - \\Delta \\left(1 + \\lambda_{2,0,0} \\kappa_{1} + \\lambda_{0,2,0} \\kappa_{2} \\right) \\right) \\right\\} ~,\n\\label{Eq:Eta3}\n\\end{eqnarray} \nwhere the prefactor in the above expressions is defined as:\n\\begin{eqnarray}\n\\Xi &=& \\Big(1 + 2 \\sum_{I} \\lambda_{I} \\kappa_{I} + 4 \\sum_{I < J} \\kappa_{I} \\kappa_{J} \\Delta_{I,J} + 8 \\kappa_{1} \\kappa_{2} \\kappa_{3} \\Delta \\Big)^2 ~.\n\\end{eqnarray}\nNote, that Eqs.\\ \\eqref{Eq:Eta1} -- \\eqref{Eq:Eta3} reduce to the two-field results in the limit $\\delta_{1,3}^{2} = \\delta_{2,3}^{2} = 0$ (cf.\\ Appendix \\ref{Sec:Scale dependent effective potential})\n\\begin{eqnarray}\n\\hspace{-20pt} \\eta_{1} &=& \\frac{16 v_{d}}{d} \\frac{ \\kappa_{2} \\lambda_{1,1}^2 + \\kappa_{1} \\left(\\lambda_{2,0} + 2 \\kappa_{2} \\Delta_{1,2} \\right)^2}{\\left(1 + 2 \\kappa_{1} \\lambda_{2,0} + 2 \\kappa_{2} \\lambda_{0,2} + 4 \\kappa_{1} \\kappa_{2} \\Delta_{1,2} \\right)^2} ~, \\\\\n\\hspace{-20pt} \\eta_{2} &=& \\frac{16 v_{d}}{d} \\frac{ \\kappa_{1} \\lambda_{1,1}^2 + \\kappa_{2} \\left(\\lambda_{0,2} + 2 \\kappa_{1} \\Delta_{1,2} \\right)^2}{\\left(1 + 2 \\kappa_{1} \\lambda_{2,0} + 2 \\kappa_{2} \\lambda_{0,2} + 4 \\kappa_{1} \\kappa_{2} \\Delta_{1,2} \\right)^2} ~,\n\\end{eqnarray}\nwhere we have written $\\lambda_{1,1} \\equiv \\lambda_{1,1,0}$, etc., as the $3$-sector effectively decouples in this case. From here, it is easy to identify the decoupling and symmetry enhancement scenarios for the spectrum of anomalous dimensions, by considering the proper limits of the couplings between different sectors. These results for the two-field case are equivalent to the anomalous dimensions given in Refs.\\ \\cite{Friederich:2010hr,Wetzel:2013}.\n\\twocolumngrid\n\n\\end{appendix}\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Approach}\n\\label{sec:approach}\nDetermining whether the bottleneck router in a network path is employing a drop-tail or AQM queuing scheme is fundamentally a classification problem, which has been rigorously studied in Machine Learning literature. To this end, we employ a data-driven approach\nthat entails generating random network configurations of varying complexity, simulating the generated scenarios using Mininet \\cite{mininet}, and training a classification algorithm on the obtained data. Our approach contains three main components: (i) randomized network topology generation and simulation, (ii) feature engineering to encode relevant information, and (iii) training of the classification algorithm. \n\n\\begin{figure}[!htb]\n \\centering\n \\includegraphics[width=\\linewidth]{figures\/topology}\n \\centering\n \\includegraphics[width=\\linewidth]{figures\/legend}\n\\caption{A random network topology generated by our framework. Edge weights denote link capacities.}\n\\label{fig:topology}\n\\end{figure}\n\n\\subsection{Random Topology Generation}\nProperly training a machine learning algorithm to simultaneously achieve high prediction accuracy and generalization to different kinds of network configurations requires a large and diverse set of training data. Hence, we opted to leverage randomization in order to generate a wide array of sufficiently diverse network topologies, which we subsequently simulated using Mininet. In particular, we developed a framework for constructing varied network configurations. \n\nAn example of a random network topology $G = (V, E)$ generated by our framework is depicted in Fig.~\\ref{fig:topology} and is constructed as follows. We initialize a \\textit{path graph}, i.e., a graph defined by a path from the source node $s$ to the sink node $g$, $\\pi = (s, v_1, \\ldots, v_n, g)$, consisting of $n \\sim \\mathcal{U}(3,5)$ \\footnote{$\\mathrm{U}(a,b)$ denotes a continuous uniform distribution, whereas $\\mathcal{U}(a,b)$ denotes a discrete uniform distribution.} switches, i.e., routers. For each switch along the path $v \\in \\{v_1,\\ldots,v_n\\}$, a random Erd\\\"{o}s-R\\'{e}nyi graph $G_\\text{rand}(n_p, p)$ \\cite{erdos} characterized by the uniformly drawn number of nodes $n_p \\sim \\mathcal{U}(1,5)$ and connectivity probability $p \\sim \\mathrm{U}(0,1)$ is generated such that each node is a switch or a host with equal probability; the largest connected component $G_\\text{conn} \\subseteq G_\\text{rand}$ is identified and connected to switch $v$ by an edge (i.e., network link). Additional $h \\sim \\mathcal{U}(1,5)$ hosts are generated and connected to vertex $v$.\n\nRandomly chosen values for link delays and capacities are generated to diversify the network configuration. For each link $e \\in E$, we set a random link delay $d \\in \\mathcal{U}(10,100)$ and link capacity $c \\in \\mathcal{U}(10,1000)$. To increase network complexity and simulate scenarios with bottlenecks caused by competing flows, we add auxiliary flows, i.e., paths, from randomly chosen host nodes to the sink node. The randomly generated link capacities are adjusted accordingly so that there exists a single bottleneck router along the path $\\pi$ from source to sink, which is subsequently labeled as either Tail Drop or AQM (denoting the queuing scheme to be simulated).\\footnote{The ranges used for the uniform distributions can be easily altered to generate random networks of desired complexity.}\n\n\\subsection{Feature Engineering}\n\n\\begin{figure}[t!]\n\\centering\n\\subfigure[Tail Drop]{%\n \\includegraphics[width=\\linewidth]{figures\/Basic-2.pdf}%\n\\label{fig:tail-drop}}\n\\subfigure[PIE]{%\n \\includegraphics[width=\\linewidth]{figures\/PIE-2.pdf}%\n\\label{fig:pie}}\n\\caption{RTT data obtained from two distinct Mininet simulations of the same network topology with the exception of the queuing scheme used by the bottleneck router: (a) Tail Drop and (b) PIE (AQM).}\n\\label{fig:data}\n\\end{figure}\nThe randomly generated network topologies are simulated using Mininet to obtain training data in the form of RTT and CWND information over time, as further detailed in Sect.~\\ref{subsect:training}. Figs.~\\ref{fig:tail-drop} and \\ref{fig:pie} depict the RTT data gathered from a 20 second Mininet simulation of the topology shown in Fig.~\\ref{fig:topology}. In order to construct a vector of distinguishing features, we examined a multitude of diversely-ranged RTT and CWND figures, similar to the figures shown above. Our observations indicated that high amount of noise was present in the data, therefore, we used the Exponentially Weighted Moving Average (EWMA) filter with parameter $\\alpha=0.3$ to smoothen the resulting RTT and CWND data.\n\nOur observations reaffirmed the inherent differences in the methodology of Tail Drop and AQM schemes: Tail Drop does not drop packets until the queue is full, which leads to a pronounced sawtooth pattern in the RTT graph (Fig.~\\ref{fig:tail-drop}), i.e., lengthy periods of linear increase in the RTT plot followed by a sharp drop. PIE, on the other hand, preemptively (and probabilistically) drops packets, which results in relatively short bursts of increasing and decreasing RTT values. \n\nWe combined our insights from experimental observations and the inherent differences in Tail Drop and AQM schemes to construct a vector of distinguishing features. Our feature vector was 72 dimensional (36 features each for RTT and CWND) that were computed for each simulated time series. Examples of particularly distinguishing features that constituted our feature vector include: statistics regarding the first and second order gradient (e.g., mean, variance, maximum), the number of local minima \\& maxima, the $L_2$-norm squared total variation, the maximum sum subsequence, and the mean and variances of the local maxima. Fig.~\\ref{fig:features} shows a complete list of the features that were used, ranked by descending importance.\n\n\\subsection{Training}\n\\label{subsect:training}\nTo train our classifier, we generated 1,100 different network topologies of varying complexity and heterogeneity so that balanced, labeled data relevant to both Tail Drop and AQM schemes were obtained. For each random network topology, we instantiated two consecutive, distinct Mininet simulations by varying the queue management technique on the bottleneck between Tail Drop and PIE. We simulated a long running TCP flow on each network from our known source to the sink and auxiliary flows from other randomly generated hosts to their respective sinks for 20 seconds. From the simulation, we gathered the resulting congestion window and RTT over time using \\verb|tcpprobe| and \\verb|ping| respectively. This data, as exemplified by the RTT plot in Fig.~\\ref{fig:data}, were then used as labeled training data in our classifier.\n\\section{Conclusion}\n\\label{sec:conclusion}\nWe presented a data-driven approach that leveraged a Neural Network for reliably determining whether a bottleneck router along a given path employs a Tail Drop or AQM scheme. We developed a framework for generating random network topologies and used Mininet to simulate a wide range of scenarios and obtain diverse training data. We identified distinguishing features by analyzing the intrinsic differences between AQM and Tail Drop and quantified the respective importance of each feature. Our feature importance analysis indicated that the most salient features were statistics pertaining to the gradient and RTT time-series data.\n\nOur algorithm attained an accuracy of 97\\% when evaluated against test data and generalized and performed well against data from highly complex topologies that it had not been exposed to before. We conjecture that our algorithm has potential to generalize to and provide reliable and accurate predictions in complex, real-world networks. Future work includes further training and evaluation on even more complex topologies and real-world data. We also plan to extend our method for classification of different AQM interfaces such as RED and PIE, in addition to the Tail Drop scheme.\n\\section{Contributions}\n\\emph{Cenk Baykal} developed the algorithm used to generate the random network configurations and helped develop the features that were used to classify the queueing scheme. \n\\\\\n\n\\emph{Wilko Schwarting} worked on the classification algorithm and evaluated state-of-the-art classification methods and conducted hyper-parameter optimization.\n\\\\\n\n\\emph{Alex Wallar} developed the Mininet codebase that was used to simulate the randomly generated network configurations and obtain data.\n\n\\section{Evaluation}\n\\label{sec:evaluation}\n\nUsing the the data obtained from our simulations, we trained and benchmarked state-of-the-art classification algorithms such as Support Vector Machines (SVMs) \\cite{svm}, Neural Networks (NNs) \\cite{hornik1989multilayer}, Logistic Regression \\cite{logistic}, and Random Forests, among others. We conducted k-fold cross validation (with $k = 10$) for each classifier on multiple data sets and found Neural Networks to be consistently outperforming all other tested classification techniques.\n\nHaving converged on using NNs as our classification algorithm, we employed a randomized search to optimize hyper-parameters by drawing from a previously specified parameter space distribution.\nThe parameter space included 1) NN architecture, i.e., the number of layers (1-4) and neurons (2-20) per layer,\n2) $L_2$-Regularization parameter $\\alpha$ penalizing large connection weights, and\n3) solvers, including Stochastic Gradient Descent (SGD), ADAM, and Limited-memory Broyden Fletcher Goldfarb Shanno (LBFGS) algorithms.\nThe randomized search evaluates each parameter tuple with a score consisting of the average accuracy over 100 stratified cross-validations with $90\\%$ training data and $10\\%$ test data, to avoid biased parameter selection from overfitting. After 10,000 evaluations, the optimal hyper-parameters were found to be a NN topology with 1 hidden layer consisting of 14 neurons, a regularization of $\\alpha = 0.5e-10$, and the LBFGS solver.\n\n\\subsection{Training Results}\nFrom the 2,200 labeled data points that were generated, 200 were left out of the training to be used later for testing. After training the classifier and optimizing its hyper-parameters, we tested how well the classification generalized by evaluating its performance on the left-out data set. Fig.~\\ref{fig:cm-left-out} shows a confusion matrix summarizing the results of this classification test. On a per-category basis, we achieve 94\\% and 99\\% prediction accuracy for correctly classifying the bottleneck router as Tail Drop or PIE respectively. Overall, we were able to accurately classify whether the bottleneck router was utilizing a Tail Drop or PIE as its queuing scheme 97\\% of the time. This can be seen as an improvement over current state-of-the-art AQM detection algorithms, such as the recently proposed method by Bideh et al. \\cite{tada}, which achieved an overall classification accuracy of 73\\% on a serial bottleneck within a relatively simple network topology. In addition to achieving a higher classification accuracy than that of previous work, we also note that our algorithm was evaluated on significantly more complex and diverse network configurations, whereas the algorithm presented in \\cite{tada} was evaluated against a single topology consisting of solely 3 hosts and 3 routers in total.\n\nTo thoroughly assess the dexterity of our algorithm in generating reliable predictions, we evaluated the classification accuracy of our algorithm on a data-set obtained from simulations of highly complex topologies with significantly more hosts, routers, and flows than the topologies that were used to train our algorithm. We emphasize the fact that our classifier was \\emph{not} trained on any data from these more complicated scenarios.\n\nDespite the challenges of this classification task, our classifier produced accurate labels for Tail Drop and PIE 80\\% and 65\\% of the time respectively, for an average classification accuracy of 73\\% overall. The confusion matrix for these tests is shown in Fig.~\\ref{fig:cm-complex}. This evaluation is evidence for our classifier's ability in generalizing well even to data obtained from highly complex network configurations that it was not explicitly trained to handle beforehand. Our favorable results demonstrate our algorithm's potential to be applied to real-world settings and achieve high prediction accuracy even against highly complex network configurations.\n\n\\begin{figure}[t!]\n\\centering\n\\subfigure[Left Out Data]{%\n\\includegraphics[width=0.49\\columnwidth]{figures\/cm_left_out}%\n\\label{fig:cm-left-out}}\n\\subfigure[Complex Topologies]{%\n\\includegraphics[width=0.49\\columnwidth]{figures\/cm_complex}%\n\\label{fig:cm-complex}}\n\\caption{Confusion matrices for the two testing scenarios, (a) data left out of the training with the same topology complexity and (b) data generated using more complex topologies than used for training}\n\\label{fig:cm}\n\\end{figure}\n\n\\subsection{Feature Importance}\nAs introduced in Sec.~\\ref{sec:approach}, our feature vector for a single data point was 72 dimensional, with 36 features each for RTT and CWND. The 72 distinguishing features for each data point mentioned in Sec.~\\ref{sec:approach} were motivated by our knowledge of the profound differences between Tail Drop and AQM schemes, as we covered in class \\cite{class}. Fig.~\\ref{fig:features} depicts a quantification of feature importance for each of the 72 features. By inspecting the highest ranking features, we obtain the insight that features pertaining to gradient-like statistics were the most distinguishing features and that data about the RTT over time were more helpful than the CWND data.\n\n\\begin{figure}[t!]\n\\centering\n\\includegraphics[width=\\columnwidth]{figures\/feature_importances}\n\\caption{Plot showing each feature's importance for classification. Note the high importance scores of features pertaining to RTT information and gradients' statistics.}\n\\label{fig:features}\n\\end{figure}\n\n\n\\section{Problem Definition}\n\\label{sec:problem}\nActive Queue Management (AQM) schemes, such as PIE \\cite{pie} and RED \\cite{red}, have been particularly effective in combating bufferbloat and attenuating network congestion. Despite their widespread success and improved performance over traditional Tail Drop queuing schemes, \nthe extent of AQM deployment in contemporary networks is not known \\cite{tada}. Detecting the presence of AQMs in a network has potential to not only provide insight into a network's internal characteristics (i.e., network tomography), but also facilitate network optimization \\footnote{e.g., in the case that a transport protocol is known to perform poorly over an AQM or certain subset of AQMs \\cite{tada}.}.\n\nIn this paper, we address the problem of determining whether a bottleneck router on a given network path is using an AQM or a drop-tail scheme. We assume that we are given a source-to-sink path of interest -along which a bottleneck router exists- and data regarding the Round-Trip Times (RTT) and Congestion Window (CWND) sizes with respect to this flow. We develop a reliable classification algorithm that solely uses RTT and CWND information pertaining to a single flow to classify the queuing scheme, Tail Drop or AQM, used by the bottleneck router. We evaluate our method and present results that demonstrate our algorithm's highly accurate classification ability across a wide array of complex network topologies and configurations.","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\t\\label{S:11}\n\t\n\tChemometric classification methods are often used to discriminate high-dimensional chemical signatures of unknown samples to determine their most likely class label \\cite{patternrec}. These methods have proven to be valuable for the fields of archeaometry and forensics where the origin, or provenance, of a manufactured item is often of interest\\cite{patternrec, obsidian, crypto, shipwreck, butter, heroin} because the chemical data, especially those obtained from neutron activation analysis (NAA) \\cite{sayre,NAArev}, tend to be highly multivariate. The high dimensionality of data poses a challenge for understanding sample relationships because it cannot be easily visualized or interpreted directly. While, these methods can succeed at classifying samples, they do so by including information from the class labels into the model, and they provide little insight about trends or patterns found in the data with the absence of labels. Information pertaining to which samples are most similar to one another between classes, or which are different from the rest of the samples within a given source class, without the influence of label discrimination is not available from classification methods. \n\t\n\tExploratory methods are often used to reveal trends and other patterns hidden in data without the use of class labels \\cite{patternrec}. Exploratory data analysis (EDA) is a form of data analysis which encompasses a number of visual exploration methods. Success in an EDA study depends on the creativity of the analyst as much as on the technique. Although there is no strict definition of EDA, it has been stated that, ``Exploratory data analysis isolates patterns and features of the data and reveals these forcefully to the analyst\" \\cite{tukeybook}, and also that ``Exploratory data analysis' is an attitude, a state of flexibility, a willingness to look for those things that we believe are not there, as well as those we believe to be there\"\\cite{tukeybooktwo}. The coupling of EDA methods with preprocessing and\/or feature selection can aid an analyst in the discovery of patterns, or the lack thereof, in data. The information gained from EDA may be used with feature selection, preprocessing and modeling methods to iteratively improve a data analysis pipeline.\n\t\n\tThe two predominant types of EDA focus on dimension reduction and clustering. These methods have been extensively used for the EDA of archeometric data \\cite{crypto, shipwreck, butter, cluster} because high dimensional data are not readily visualized. Dimension reduction techniques such as principal component analysis (PCA), discriminant analysis, or exploratory projection pursuit are used to change the basis of the data to one based on lower-dimensional projections \\cite{patternrec}. Unfortunately, these methods do not offer a guarantee that the resulting 2-D or 3-D projections will offer meaningful information about class relationships or sample-class relationships, due to the presence of class overlap or class heterogeneity\\cite{patternrec}. Unsupervised classification as implemented in clustering algorithms suffers from the opposite problem clustering methods readily create groupings of high-dimensional data but tend to offer little information about the cluster assignments themselves. Clustering methods have been known to incorrectly associate archeological samples depending on the method used and on the chemistry of the samples \\cite{cluster}. Information about sample or class relationships may be inferred by employing many techniques, but the resulting information is often hidden, and typically key relationships in the data cannot be seen by an analyst \\cite{cluster}.\n\t\n\tThis work reports a new approach, which we call the uncharted forest, to the visualization and measuring of relationships within and between classes of data. This approach has elements of a clustering algorithm in that it groups similar samples with one another, but is also similar to dimension reduction methods in that it outputs a single heat map which can be interpreted to reveal information about the samples. Uncharted forest analysis uses a partitioning method that is related to the sample partitioning approach used in decision trees but, it does not use class labels like most tree methods do \\cite{CARTBook}. Instead, the uncharted forest analysis explores how samples relate to one another under the context of univariate variance partitions. Although the method is unsupervised, we show that when the results are overlaid with external class labels, the method can be used to investigate sample relationships as they pertain to class labels. We demonstrate that this technique can be used as a tool for exploratory data analysis to visualize class or cluster associations, sample-sample associations, class heterogeneity, and uninformative classes. The utility of uncharted forest analysis is demonstrated on two classification datasets and two provenance datasets. Additionally, two empirical clustering metrics are compared with two of the metrics obtained by uncharted forest analysis on another provenance dataset.\n\t\n\t\n\t\\section{Theory}\n\tTo motivate the use of an unsupervised tree ensemble, a brief review of supervised trees and bagged tree ensembles as used in classification is provided here. A complete review of these methods can be found elsewhere \\cite{CARTBook,BreimanRF}. \n\t\n\t\\subsection{Supervised Classification Trees}\n\tBefore describing the mechanisms by which supervised decision trees are used to partition data, a survey of the vocabulary common to these methods is presented. Classification trees are a type of supervised classifier, which means that the assignment of new data to a class label requires that each sample the tree was trained on also has a label. The aim of developing a classification tree is the establishment of a set of sequential rules that can be applied to label a sample based on its features. Supervised decision trees are collections of many binary decisions, where each decision is made at one of three locations: roots, branches, or terminal nodes\/leaves. These three locations are displayed in their hierarchical ordering in Figure 1. A root is the first decision made in the tree. Branch nodes indicate later, non-terminal decisions. Terminal nodes indicate where a branch has been terminated and where final decisions that assign class labels are made. \n\t\n\t\\begin{figure}[h]\n\t\t\\centering\n\t\t\\includegraphics[width=0.40\\linewidth]{Figure1}\n\t\t\\caption{A block diagram showing the relationships between roots, branches and terminal nodes via arrows.}\n\t\\end{figure}\n\t\n\tEvery sample from the training set is partitioned at a branch or node based on whether a selected variable for the sample has a value that is greater than or less than a specified value. The combination of a variable and value that partitions a set of samples is often referred to as a decision boundary. Decision boundaries are obtained by exhaustively searching each variable and finding the threshold value which affords the highest gain \\cite{CARTBook}. Here, gain is a metric-dependent measure of how well a decision separates the available samples according to their class labels. Metrics that can be used to minimize class label impurities at terminal nodes, such as the Gini impurity, informational entropy, and classification accuracy are commonly used to train supervised classification trees by finding a set of decision boundaries \\cite{CARTBook}. \n\t\n\tFor example, if the samples available at a branch node are better separated into their respective classes by application of some decision, two new branches are made at that node. If not, the partitioning terminates, the node becomes a terminal node, and the decision tree ceases growth using those samples. In this way, each observation is sorted from the root of a tree through branches and eventually to terminal nodes, based on whether creating a new decision boundary better separates samples by known class labels\\cite{CARTBook}. \t\n\t\n\tOne advantage to the use of decision trees is that the classifier that results can be easily understood by examining the decision boundaries that partition the data and the relation of these boundaries to class labels. However, the predictive performance of classification trees on new samples is generally poor relative to other models because the partitions are sensitive to noise in the variables \\cite{CARTez}. Tree-based classifiers also tend to overfit the data, and some sort of branch removal, or pruning step, is usually employed to reduce the complexity of the tree \\cite{CARTBook}. \n\t\n\t\\subsection{Random Forest Classification}\n\tThe random forest classifer is a multiclass classification technique that utilizes the results from an ensemble, or collection of supervised decision trees \\cite{BreimanRF} to assign class labels. In this technique, a user-specified number of supervised decision trees is created, each overfit to a selection of bootstrapped samples and random selections of variables. Random forests make use of the complexity of many, biased decision trees to make robust classification assignments by using the average class designation assigned by every tree to predict overall class membership for new data. This approach to ensemble model building is referred to as bootstrap aggregating, or bagging \\cite{bagging}. The average of many, complex trees tends to perform well in practice \\cite{BreimanRF}.\n\t\n\tRandom forest classification has some distinct advantages relative to many other classifiers. It is robust to outliers because sample selection at branches is bootstrapped, and the effect of a few outlying samples does not affect its cost function significantly \\cite{BreimanRF}. Random forest models are also robust to noise in variables because many classifiers, each with randomly selected variables, are used to decide on class memberships. Another notable advantage to random forest classification is that there is no requirement for linearity in class boundaries. Random forest models, unlike methods such as linear discriminant analysis or support vector machines, are not based on the hypothesis that classes are linearly separable \\cite{linearmethods,kernel}. Random forest classification also requires relatively few hyperparameters and works without the need for significant tuning of the classifier on many kinds of data \\cite{CARTez}. \n\t\n\t\\subsection{Uncharted Forest Tree}\n\tThe concept of unsupervised decision tree modeling is relatively new to chemometric applications. Only a few approaches using unsupervised trees have been reported. An approach reported by Khudanpur, et al. is based on an algorithm similar to the one reported here, but their approach uses the Kullback-Leibler distance to find sets of unsupervised partitions that are optimal with respect to an information theoretic measure\n\t\\cite{InformationApp}. Other methods for developing unsupervised decision trees are based on the use of an assumed class label for each observation; the trees are trained in a manner that is similar to that used for training a random forest classifier \\cite{UnsupRF, CLTrees}. \n\t\n\tOur approach to unsupervised decision trees, the uncharted forest tree algorithm, focuses on intuitive concepts from statistics or machine learning rather than on information theory, to allow for easy interpretation. The trees can be created without the need for class labels because the gain function that is optimized in the construction of each tree relies only on information in the data matrix, not on the labels. This algorithm does not utilize the usual supervised metrics for optimization such as the Gini importance or entropy \\cite{BreimanRF}. Instead, the tree hierarchies used in uncharted forest are constructed from decision boundaries that reduce measures of spread in a given variable. The reduction of spread is common in pattern recognition and chemometrics; it is the mechanism underlying techniques such as K-means clustering, and principal component analysis for dimension reduction \\cite{patternrec}. Metrics based on changes in the variance, median absolute deviation (MAD), or mean absolute deviation (AD) can also be used for creating the decision boundaries used in an uncharted forest tree. If variance is used, however, univariate Gaussian probabilities may be calculated from the samples on either side of a decision boundary. We prefer the variance metric because it is related to several established figures of merit for clustering, as we describe later. The gain metric based on changes in variance that we use is\n\t\n\n\t\\begin{equation}\n\tgain = 1 - \\frac{Var(Branch_1) - Var(Branch_2)} {Var(Parent Node)}\n\t\\end{equation}\n\t\n\t\\noindent where the sample data in either partition are denoted as $Branch_1$ and $Branch_2$, and the sample data prior to introducing the partition is contained in the previous - or parent - node, which has been denoted as, $Parent Node$. The gain in the unsupervised case can then be seen as the loss of variance ($Var$) in the data resulting from the construction of a given decision boundary. The gain in the unsupervised case can then be seen as the loss of variance in the data resulting from the construction of a given decision boundary, relative to that in the previous - or parent - node. Similar gain equations can be written for measures of spread based on AD and MAD. These are not considered in this work because they do not have direct comparisons to clustering metrics. \n\t\n\tThe spread in data can almost always be reduced by partitioning data around an available decision boundary. Thus, uncharted forest trees require a user-defined parameter, the tree depth, to limit the number of possible partitioning decisions made. Tree depth can be defined as the maximum number of branches that are counted from the root node to a given terminal node. For exploratory analysis of classification or cluster data, a reasonable depth may be obtained from the rounded base 2 logarithm of either the number of known classes or the number of clusters selected. The number of classes\/clusters available is often known information; for example, in most provenance studies the number of classes is decided upon by the number of sampling sites used. In our implementation the minimum node size, or when branch growth is ceased after the number of samples becomes fewer than a given number, has been set to five samples because five samples allows for a calculation of variance that is not undersized but still accommodates a small cluster\/class size. Pseudocode that depicts the recursive algorithm for constructing an uncharted decision tree is shown in Algorithm 1.\n\t\n\t\n\n\t\n\t\\begin{algorithm}\n\t\t\\caption{Uncharted Forest Decision Tree}\\label{unchartedtree}\n\t\t\\begin{algorithmic}[1]\n\t\t\t\\Procedure{ }{}\n\t\t\t\\State for each variable in variables\n\t\t\t\\State \\ \\ \\ \\ Calculate gain. (Eqn. 1)\n\t\t\t\\State \\ \\ \\ \\ Store highest gain.\n\t\t\t\\State end for\n\t\t\t\\State Partition samples at largest gain.\n\t\t\t\\If {branch depth $\\geq$ number of branches} \n\t\t\t\\State Stop growth.\n\t\t\t\\Else\n\t\t\t\\State For each partition \\textbf{goto} 2\n\t\t\t\\EndIf\n\t\t\t\\State \\textbf{end if}\n\t\t\t\\EndProcedure\n\t\t\\end{algorithmic}\n\t\\end{algorithm}\n\t\n\tThe uncharted forest decision tree partitions a set of data by finding decision boundaries that reduce the spread in the data. It functions as a kind of clustering, where the clusters result from divisive partitioning of spread in the data. Unfortunately, unsupervised trees of the form described above often have marginal utility in practice. These trees can be used for clustering but, like a single supervised tree, the individual uncharted forest decision trees are sensitive to noise effects in the data. The variables in a set of data have many contributions to their spreads, and as a consequence, the partitions that result from a single unsupervised tree tend to be uninformative.\n\t\n\t\\subsection{Uncharted Forest}\n\tOne solution to the uninformative partitions created by a single uncharted forest decision tree is to create an ensemble of unsupervised trees, an ``uncharted forest\". The approach is similar to that taken in the random forest classifier. For each tree in this uncharted forest, the variables used are randomly selected, so that potentially relevant information within an individual variable is not obscured by other contributions to its spread or by another variable with much larger gain in the reduction of spread. The algorithm for the uncharted forest method follows that of the random forest algorithm \\cite{BreimanRF}, except that the supervised classification and regression tree models are replaced with the uncharted forest decision tree and the samples are not bootstrapped.\n\t\n\tBecause construction of the trees is unsupervised, the output of the uncharted ensemble algorithm differs from that of the random forest. An uncharted forest returns a square matrix whose dimensionality is that of the number of observations, so that each row or column vector stores information about a single observation and how it relates to the other samples. The matrix is populated in the following manner: when samples are partitioned into a terminal node together, the row and column representative of each of these samples each gain a +1 count to account for their associations with one another. Similarly the row and column of any samples that were not present at a terminal node with those samples each receive a -1 count. The associations formed by random chance become less pronounced by including the -1 penalty on observations that were not grouped with one another at their respective terminal nodes. After row-wise normalization, the matrix possesses diagonal symmetry. In the case of variance-only partitions, a Gaussian probability for node membership may be used instead of $\\pm$ 1. \n\t\n\n\t\n\t\\subsection{Figures of Merit}\n\tWe refer to the matrix generated by the tree ensemble as a probabilistic affinity matrix, because of its resemblance to the distance-based affinity matrices utilized by clustering methods \\cite{spectral}. Each element in the matrix represents the row-wise, likelihood that a given row (observation) is associated with a given column (another or the same observation) at a terminal node in an uncharted forest tree. In order to maintain the above interpretation, a ramp function is applied to every matrix element (p), max(0, p), to remove any negative entries. In this way, the diagonal elements have a sampled probability of association of 1.0, because there is a 100\\% likelihood that a given sample will be associated with itself at each terminal node where it was assigned. \n\t\n\tWe introduce two schemes that allow for the probability values in the matrix to be quantified for exploratory analysis. Two figures of merit can be calculated to measure how often samples between- or within- clusters\/classes are partitioned together. The submatrix of the probability affinity matrix (P) that contains the associations between two clusters\/classes i and j, can be quantified by what we define as the interaction quotient (IQ),\n\t\n\n\t\\begin{equation}\n\tIQ_{i,j} = \\frac{\\sum_{c=j_{min}}^{j_{max}} \\sum_{r=i_{min}}^{i_{max}} P_{r,c}} {(j_{max} - j_{min}))(i_{max} - i_{min}) }\\end{equation}\n\t\n\t\\noindent where the min and max subscripts denote the contigious row\/column beginning and end for cluster\/class indices within the matrix. The IQ is simply the sum of all probabilities inside of the submatrix that is defined by the rows shared from cluster\/class $i$, and the columns shared by cluster\/classes $j$ normalized by the maximum possible probabilities for association within that submatrix. Similarly, we define the self association quotient (SAQ) as a measure of within- class or cluster sample associations. The SAQ metric is calculated from the probabilistic affinity matrix for the special case where $i = j$. For purposes of visualization, the matrix blocks from which these metrics are calculated are displayed below,\n\t\n\t\\[\n\t\\begin{bmatrix}\n\tSAQ_{1} & IQ_{1,2} & IQ_{1,3} & \\dots & IQ_{1,j} \\\\\n\tIQ_{2,1} & SAQ_{2} & IQ_{2,3} & \\dots & IQ_{2,j} \\\\\n\t\\vdots & \\vdots & \\vdots & \\ddots & \\vdots \\\\\n\tIQ_{i,1} & IQ_{i,2} & IQ_{i,3} & \\dots & SAQ_{i}\n\t\\end{bmatrix}\n\t\\]\n\t\n\tThe total interaction quotient (TIQ) which describes the extent of the total interactions between the samples of all class\/cluster assignments can be expressed by the sum of probabilities within each IQ block divided by the number of matrix entries in those blocks(n),\n\t\n\n\t{\\centering\n\t\t\\begin{equation} TIQ = \\sum_{i} \\sum_{j} IQ_{i,j} \/ n \\ \\ \\ \\ \\ provided \\ \\ i \\neq j \\end{equation} \\par\n\t}\n\t\n\tIn an analogous manner, the total self association quotient (TSAQ) can be used to describe how often the samples in all class\/clusters are partitioned with samples in their same respective class\/cluster by uncharted forest,\n\t\n\n\t\\begin{equation}TSAQ = \\sum_{i} SAQ_{i} \/ n \\end{equation}\n\t\n\tThe metrics SAQ and IQ provide a means to assess the likelihoods that samples are within the same terminal nodes after variance partitioning with reference to class\/cluster ownership. SAQ can provide a measure of class\/cluster heterogeneity by indicating how likely it is that the samples inside of a given label are partitioned together. The IQ provides a measure of how often data were associated to the samples that belong to a different class\/cluster label. The IQ measure can therefore be used as an aid to identify classes in the data that do not associate with samples from other classes of data, or classes which are similar to one another. These measures can also be made on a sample-wise basis to discover outliers which may either be between class\/cluster labels or are not associated with any class\/cluster label at all.\n\t\n\t\n\t\\section{Scope of the Uncharted Forest Method}\n\tThere is one restriction to the use of uncharted forest analysis that should be noted. The associations that a uncharted forest analysis can describe are often not meaningful without data ordered by class or cluster labels, despite that it is an unsupervised technique. The hypothesized membership of samples allow interpretations of the data as it pertains to class\/cluster labels. In some cases, data have serial correlations, such as in time series, or variables from spectroscopic instruments, and the probabilistic affinity matrix appears well-ordered. However, with uncorrelated, labelless data, the probabilistic affinity matrix most often appears similar to random noise. One interesting feature of the probabilistic affinity matrix is that it is invariant to the ordering of data. This is a useful attribute when the matrix could display seemingly random disjoint probabilities, because the order of the rows and columns of the matrix can be changed to permit a more meaningful interpretation, like a cluster for example, without running the algorithm again. The studies explored in this work address the case where class or cluster labels were available for provenance and classification data.\n\t\n\t\\section{Selected Examples}\n\tIn this work, five datasets were investigated by uncharted forest analysis. The first, the well-known iris dataset, is presented as an instructional example to demonstrate how the technique can be used to examine data with class labels, and how it compares to the widely used principal components analysis method. The well-known ARCH obsidian data is presented to compare the sample-sample associations obtained by uncharted forest analysis with two supervised classification methods. The Nan'ao One shipwreck and cryptocrystalline sillicate datasets are presented to compare the results from uncharted forest analysis to those of published works that used other data analysis methods on the same data. The fifth dataset, the butter data, is used to compare the new method to established variance-based partitioning methods and to show that the TIQ and TSAQ figures of merit are similar to well-established, variance-based clustering metrics.\n\t\n\t\\subsection{Iris Data EDA}\n\tThe iris dataset contains measurements of 150 iris flower samples. Both the widths and lengths of the sepals and petals of each flower were measured. Each sample belongs to either the virginica, versicolor, or setosa species \\cite{iris1, iris2}. The iris data is widely regarded as a standard dataset for introducing classification problems and algorithms because it features two elements common to classification problems: it possesses a linearly separable class (setosa) and two classes with overlapping boundaries (virginica, versicolor).\n\t\n\tWhen the uncharted forest algorithm was applied to the iris data using 100 trees, a depth of 4, and using the variance metric, the probabilistic affinity matrix indicated associations between several samples of virginica and versicolor. The row-wise total interaction quotients for the samples in the setosa class were $\\ll$ 5\\%; this suggested that the setosa class was well separated from the others. Four samples that were greater than 3$\\sigma$ from the mean row-wise IQ were found in both the virginica and versicolor classes. This result implies that, by making consecutive decision boundaries to reduce variance, the virginica and versicolor classes had several samples that appeared to share the same class label space. \n\t\n\tPrincipal components analysis was employed to reduce the dimensions of the iris data and to visually assess which samples uncharted forest analysis indicated as associated with other classes. When the iris data were projected onto the first two principal component axes, five samples which had significant row-wise IQ's were located between the virginica and versicolor classes, as can be seen in Figure 2. A conventional random forest model with 500 trees was then trained on the entire dataset. The random forest model misclassified 5 samples. Of those 5 samples, 3 were the samples with row-wise IQ values that exceeded 3$\\sigma$. The heat map from uncharted forest algorithm highlighted many of the samples that the random forest classifier had difficulty classifying due to class overlap.\n\t\n\t\\begin{figure}[H]\n\t\t\\centering\n\t\t\\includegraphics[width=0.9\\linewidth]{Figure2}\n\t\t\\label{fig:Ng1} \n\t\t\n\t\t\\caption{A heat map of the probabilistic affinity matrix of the Iris data with class labels (top left). The row-wise IQ values for the corresponding heat-map with the 3$\\sigma$ vertical line (top right). Plot of the PCA scores for the iris data. Samples denoted with an X, are the samples with row-wise IQ values above the 3$\\sigma$ line (bottom).}\n\t\\end{figure}\n\t\n\t\\subsection{Obsidian Data Pattern Recognition}\n\tThis dataset contains X-ray fluorescence measurements made on 12 obsidian artifacts and 63 obsidian samples that were collected from four locations in Northern California. The four source locations were from quarries located near: Mount Konocti, Borax Lake, Glass Mountain (Napa County), and Anadel (Sonoma County). We used this dataset demonstrate that the sample-sample associations afforded by uncharted forest analysis could be used to establish the likely origin of the artifacts from measured quantities of the following trace elements: iron, titanium, barium, calcium, potassium, manganese, rubidium, strontium, yttrium, and zirconium \\cite{Obsidian}. This dataset has previously been demonstrated as a test-bed for unsupervised and supervised learning methods in chemometrics \\cite{kowalski}.\n\t\n\tThe unknown samples were assigned source labels based on an uncharted forest model made with 200 trees and a tree depth of 2, as well as a nearest neighbor classification (k = 1) after PCA projection onto the first two components of the unknown samples The assignments of the artifact samples to their predicted source classes can be seen in Figure 3. The PCA was created from standard normal variate-scaled obsidian source data that had the X-ray fluorescence measurements of iron, strontium, and manganese removed due to correlation with other variables. The uncharted forest analysis assignments were made from the row-wise IQ values of the rows in the probabalistic affinity matrix, which corresponded to the twelve unknown samples via a maximum voting scheme. The voting scheme was performed to remove analyst bias and to automate the assignment of samples to suspected classes, but the results were equivalent to our visual assessment. \n\t\n\t\\begin{figure}[H]\n\t\t\\centering\n\t\t\\includegraphics[width=0.95\\linewidth]{Figure3}\n\t\t\n\t\t\\caption{The probability affinity matrix (left) and the PCA scores plot obtained from the first two principal components of the standard normal variate-scaled ARCH data (right). The unknown artifact samples are denoted with a `?', and the four source classes of the obsidian data are labeled as follows: Mount Konocti (1), Borax Lake (2), Glass Mountain (3), and Anadel (4). }\n\t\\end{figure}\n\t\n\tThe source assignments of the artifacts from the uncharted forest analysis with the voting scheme differed from those made from the nearest neighbor (1-NN) PCA classification model by only one sample. For a sense of comparison, an ordinary random forest classification model with 200 trees was also trained on the source data. Again, the assignments made by the maximum vote uncharted forest analysis differed by only one sample from those of the random forest model. However, the classification produced by the nearest neighbor classifier differed from the assignments made by the random forest model by two samples. None of the models assigned any of the artifact samples to the Anadel source class. This result is consistent with a classification analysis performed by a linear learning machine \\cite{kowalski}.\n\t\n\tSimilar assignments between unknown samples and assigned source classes were obtained by the uncharted forest method and the two supervised classification models. This finding suggested that the unsupervised partitioning method used in the uncharted forest method with a maximum voting scheme can offer sample-sample associations that are comparable to that found using supervised methods.\n\t\n\t\n\t\\subsection{Nan'ao One Shipwreck Data EDA}\n\tThis dataset was selected because it was relatively high in dimension, and because it possessed complex class labels and sample membership. Twelve porcelain samples were collected from a shipwreck and compared, by using NAA measurements of 27 elements, with samples from several nearby kilns. The data were scaled by applying a base ten logarithm prior to analysis as suggested by the authors \\cite{shipwreck}. The goal of the study was to determine which kilns were most likely used to create the blue and white porcelain artifacts that accounted for 95\\% of the samples found on the wreck \\cite{shipwreck}.\n\t\n\tAn uncharted forest analysis using 30 trees and a tree depth of 2 was used to investigate the source and artifact data, as shown in Figure 4. It was found that two of the artifact samples had high associations with source samples obtained from the Jingdezhen kiln, and no associations with other source samples. This result implied that those samples could be assigned to the Jingdezhen kiln. Zhu et al.'s PCA analysis agreed with the assignment of those two samples to the Jingdezhen kiln \\cite{shipwreck}. \n\t\n\t\\begin{figure}[H]\n\t\t\\centering\n\t\t\\includegraphics[width=0.65\\linewidth]{Figure4}\n\t\t\n\t\t\\caption{Probability affinity matrix of the three source classes and the shipwreck samples. }\n\t\\end{figure}\n\t\n\tThe probabilistic affinity matrix also indicated that samples from the Jingdezhen kiln may be composed of two subclasses. Our PCA analysis did not readily show that the Jingdezhen kiln could be represented as two subclasses. It was found that the scores of the subclassess that were identified by uncharted forest analysis were linearly separable after projection from the fourth and fifth principal components. Manual assessment of the 729 pairs plots showed that the two subclasses identified by uncharted forest were linearly separable in the Lutetium measurements. The base-10 logarithmically scaled Lutetium variable was was used in 16\\% of all decision boundaries made from the random selections of variables and was the third most used variable for constructing decision boundaries in all of the uncharted forest trees. In every instance the decision boundaries were created at -0.569, and this boundary was never imposed at a root node. A plot of normalized kernel density estimates of Lutetium, displayed in Figure 5, showed that the samples from these two possible subclasses were well-separated by the Lutetium decision boundary. This information was readily accessible by uncharted forest analysis, whereas, by the other EDA methods we did not know that this division of samples in the Jingdezhen class may be of interest, or that it existed. The reason why the Jingdezhen class has two well-separated distributions of Lu measurements is not known, but this finding could be of value for future investigations or to provide information about appropriate sample selection that accounts for the sources of variation, as is necessary in the training of classification models. \n\t\n\t\\begin{figure}[H]\n\t\t\\centering\n\t\t\\includegraphics[width=0.75\\linewidth]{Figure5}\n\t\t\n\t\t\\caption{ Kernel density estimates for the base-10 logarithmically scaled Lutetium concentrations in the speculated Jingdezhen 1 and 2 subclasses as well as in the other samples. The kernel density estimates were created using a bandwidth of 0.04 and were normalized relative to the kernel estimates of the entirety of the data. The decision boundary created by the uncharted forest algorithm is shown as a vertical line. }\n\t\\end{figure}\n\t\n\tUncharted forest analysis revealed that the remaining artifact samples were found to be most associated with those from the Zhangzou province. However, there was notable interaction between samples from the Zhangzou province and three source samples from the Dapu kiln. Similarly, Zhu et al. found that the remaining samples appeared to be from the Zhangzou province, and that several samples were outside of the 90\\% confidence bounds of their principal component analysis \\cite{shipwreck}. \n\t\n\tTo make the partitions of uncharted forest more specific to the samples believed to be in the Zhangzou province, both the source and artifact samples attributed to the Jingdezhen kiln were removed. The source samples collected from the Dapu kiln were also removed to focus the partitioning only on the kilns in the Zhangzou province. The goal of this analysis was to assess the samples obtained from specific kilns within the Zhangzou province as they related to the shipwreck, rather than assuming the Zhangzou province was a homogeneous collection of source kilns. The probabilistic affinity matrix, shown in Figure 6, produced by uncharted forest analysis with 75 trees and a tree depth of 4 displayed no associations (IQ = 0) between the samples from the shipwreck and the Hauzilou and Tiankeng kilns. This finding suggests that both kilns were unlikely sources for the artifacts, a finding that was not observed in the original study\\cite{shipwreck}.\n\t\n\t\\begin{figure}[H]\n\t\t\\centering\n\t\t\\includegraphics[width=0.65\\linewidth]{Figure6}\n\t\t\n\t\t\\caption{The probabilistic affinity matrix after the fragmentation of the Zhangzou province into its respective source kilns and the removal of the two artifact samples associated with the Jingdezhen kiln, and the sources from the Jingdezhen kiln and the Dapu kiln. }\n\t\\end{figure}\n\t\n\t\n\tAlthough the additional analysis did not clearly display which individual kilns in the Zhangzou province were most likely associated with the artifact samples, it did discount two unlikely locations. In a provenance study it is also of interest to know where samples likely did not originate from, because that information can be used for the design of future studies.\n\t\n\t\\subsection{Cryptocrystalline Sillicate Data EDA}\n\tHoard et al. investigated the geological source of several cryptocrystalline sillicate artifacts that appeared identical by eye and by microscope \\cite{crypto}. Artifacts were collected in Eckles, Kansas and at three sites in Nebraska (Signal Butte I\/II\/III). Eleven trace elements in the artifacts were quantified by NAA and compared to source materials collected from Flattop Butte, Gering formation, Nelson Butte, Table Mountain, and from West Horse Creek. \n\t\n\tThe Gering formation TIQ was calculated to be $\\ll$ 5\\% by an uncharted forest model that had 100 trees and a tree depth of 4. This result suggested that the Gering formation was not commonly partitioned with samples outside of its assigned class label, in agreement with the authors' key finding that the Gering formation was not a likely source for the artifacts\\cite{crypto}. This finding is a clear example of the identification of an uninformative class by uncharted forest analysis. This finding was corroborated from the scores obtained by a PCA constructed using only the source classes. However, the PCA scores plots were uninformative with regards to obtaining other information Hoard et al. had described using canonical discriminant analysis (CDA), or that observed via the probabilistic affinity matrix obtained from uncharted forest analysis. The score plots of the first four principal components showed nondistinct overlap of almost all of the samples belonging to the classes, even after the CDA-separable Gering formation class was removed. \n\t\n\tThe probabilistic affinity matrix obtained by uncharted forest analysis using 100 trees and a tree depth of 4 with only the Eckles artifact class is displayed in Figure 7. This Figure clearly shows that the most likely source for the samples collected at the Eckles site was the Flattop Butte. This finding by uncharted forest concurred with the CDA analysis performed by Hoard et al \\cite{crypto}. \n\t\n\n\n\n\n\t\n\n\n\t\n\t\\begin{figure}[H]\n\t\t\\centering\n\t\t\\includegraphics[width=0.70\\linewidth]{Figure7}\n\t\t\n\t\t\\caption{The probabilistic affinity matrix from uncharted forest clustering of the cryptocrystalline sillicate data with only the Eckles Kansas artifact and the source classes included. Class boundaries are shown by the overlays of vertical and horizontal lines. }\n\t\\end{figure}\n\t\n\tHowever, unlike the CDA analysis performed by Hoard et al. or the PCA analysis explored in this study, an uncharted forest analysis displayed class heterogeneity. A very clear example of a class with heterogeneous structure, shown in Figure 8, was observed via uncharted forest analysis using 100 trees and a tree depth of 4, by comparing associations between artifacts from the Signal Butte III site and the source locations. Seven of the twelve samples collected from the Signal Butte III site were most frequently partitioned with two samples collected from Table Mountain, not the Signal Butte III samples. Those two samples that were collected from Table Mountain did not appear to be associated with the rest of the data that was labeled as Table Mountain. This instance of class heterogeneity may be indicative of a missing source class, the introduction of a sampling error, or of another anomaly in the Table Mountain data. Hoard et al. was unable to determine the likely origin for several of the samples that were collected from the Signal Butte III site; CDA indicated that the samples from Signal Butte III were tightly grouped \\cite{crypto}. Uncharted forest analysis, on the other hand, indicated that the samples may actually be from two groups, one not well-represented in the NAA data, and the other most like the samples obtained from Table Mountain.\n\t\n\t\\begin{figure}[H]\n\t\t\\centering\n\t\t\\includegraphics[width=0.70\\linewidth]{Figure8}\n\t\t\n\t\t\\caption{The probabilistic affinity matrix from uncharted forest clustering of the cryptocrystalline sillicate data, with overlaid class boundaries. Heterogeneity within the first 7 samples collected from the Signal Butte III site can be observed in that they only associate with two samples collected from Table Mountain. These samples do not associate with the rest of the data from Table Mountain. }\n\t\\end{figure}\n\t\n\t\n\t\\subsection{Geographical Butter Data Clustering Analysis}\n\tThe butter dataset consists of 78 butter samples collected from 17 countries. The goal of this study was to assess the viability of determining the origin of butter samples from their chemical signatures \\cite{butter}. Butter samples were analyzed by inductively coupled plasma mass spectrometry (ICP-MS) and isotope-ratio mass spectrometry using isotopic ratios of $^{13}$C, $^{15}$N, $^{18}$O, and $^{34}$S \\cite{butter}. Of the 17 countries included in the data, 8 were represented by two samples or less; these data were removed. A total of 65 samples were retained for analysis, of which 44.6\\% were of German origin. A significant proportion of the measurements of $^{34}$S in butter proteins were lower than the detection limit \\cite{butter}. To make use of the information in this variable, the samples with $^{34}$S values lower than the detection limit were imputed with zero concentrations. \n\t\n\tIn this study, the butter dataset was used to compare the variance partitions made by the uncharted forest algorithm with two variance-based clustering algorithms. K-means clustering and hierarchical clustering with a Ward's criterion were applied to the butter data, using a range of clusters (2-7), with 15 replicates. An uncharted forest analysis with 200 trees and a tree depth of 5 was used to assess each replicate of the clustering methods. K-means clustering optimizes the location of cluster partitions by minimizing within-cluster variance in an iterative fashion \\cite{kmeans}. Ward's clustering was employed with a loss function that minimizes within-cluster variance by forming clusters using an agglomerative approach\\cite{ward}. The goal of this study was to assess how the TIQ and TSAQ metrics compared to within- and between-cluster variance metrics, respectively. \n\t\n\tConfidence intervals of the correlation coefficients between uncharted forest metrics and the empirical cluster variances for both clustering methods were created by Fisher transforming the correlation coefficients, calculating the confidence intervals as usual, and transforming back \\cite{Fisher}. All intervals reported here were computed at 99\\% confidence. TSAQ and within-cluster variance were found to be strongly correlated: K-means clustering had a correlation of 0.90 $\\pm$ 0.12 with TSAQ, and Ward's clustering had a correlation of 0.94 $\\pm$ 0.05. This result strongly suggested that the relationship between TSAQ and the within-cluster variance metric were well approximated by a linear relationship. The intervals for the correlation coefficients between TIQ and between-cluster variances were -0.50 $\\pm$ 0.07 for K-means clustering and -0.60 $\\pm$ 0.05 for Ward's clustering. Overall, the correlation between TIQ and between-cluster variance of either clustering method was moderate; this finding indicated that TIQ provides information that is only somewhat similar in direction to between-cluster variance.\n\t\n\tThe negative correlation between TIQ and between cluster variances was expected, because as the number of clusters increases, the cluster centroid distances also increase, and so do the between cluster variances. However, as cluster distances increase, the associations found by uncharted forest analysis between the samples inside of those clusters would be expected to become less prevalent, because the clusters are becoming more specific to the groups in the data; thus, TIQ values would be expected to decrease. The magnitude of the correlations were not near-perfect (-1.0) because TIQ and TSAQ metrics result from randomly created partitions and the uncharted forest algorithm does not partition data the same way that either K-means or Ward clustering does. However, the fact that the correlation was correct in sign suggests that the partitioning mechanism is overall variance-based, as intended.\n\t\n\t\n\t\\section{Conclusion}\n\tWe have presented a novel adaptation to the random forest paradigm which extends it to unsupervised exploratory data analysis. The method allows for high- and low-dimensional data to be represented as a matrix where every entry is a sampled probability-like value that represents the likelihood that a given sample (row) resides in the same terminal node as the other samples (column). This exploratory data analysis tool and its associated metrics are shown to be a informative aid for the analysis of classification and provenance data. Uncharted forest analysis was able to provide information related to uninformative class labels, class heterogeneity, and sample-wise relationships. \n\t\n\tThe uncharted forest analysis method using variance partitions was explored by comparing its measurable results with both between- and within- cluster variances of two known variance based clustering algorithms. It was shown that the uncharted forest algorithm can provide metrics (TIQ, TSAQ) that are correlated to global clustering statistics. However, the true utility of the method is found in the visualization of sample-wise information, and how it relates to cluster\/class-wise associations. As a proof-of-concept though, when sample-wise variance partitions were used in uncharted forest analyses, the probabilistic affinity matrix generated metrics that were correlated to rigorously defined statistical measures of spread for clustered data. \n\t\n\tVery similar findings to the classification based provenance studies performed by Zhu et al. and Hoard et al. were obtained from exploratory uncharted forest analysis. This is of interest because the uncharted forest algorithm is unsupervised, but with an overlay of class labels, it was possible to visually infer class labels for unknown samples in high-dimensional, multiclass datasets. Most importantly, uncharted forest analysis was able to display information about the data that was not readily apparent from the classification methods used in the aforementioned studies and other exploratory data analysis methods. The uncharted forest method was shown to be an EDA technique that can provide information that is relevant to the domains of exploratory data analysis, clustering, and classification.\n\t\n\t\\section{Acknowledgments}\n\tThis work was supported by the United States National Science Foundation grant 1506853.\n\t\n\t\\section{References}\n\n\n\n\t\n\n\n\t\n\n\n\n\n\n\n\n\t\n\n\t\n\t\\bibliographystyle{model1-num-names}\n\t","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\n\\section{Introduction}\n\\label{sec:intro}\nIn this work, we study local problems on regular trees from three different perspectives. \n\n\\textit{First,} we consider the perspective of distributed algorithms. In distributed computing, the studied setup is a network of computers where each computer can only communicate with its neighbors. Roughly speaking, the question of interest in this area is which problems can be solved with only a few rounds of communication in the underlying network. \n\n\\textit{Second}, we consider the perspective of (finitary) factors of iid processes.\nIn probability, random processes model systems that appear to vary in a random manner.\nThese include Bernoulli processes, Random walks etc.\nA particular, well-studied, example is the Ising model.\n\n\\textit{Third}, we investigate the perspective of descriptive combinatorics. The goal of this area is to understand which constructions on infinite graphs can be performed without using the so-called axiom of choice.\n\nAlthough many of the questions of interest asked in these three areas are quite similar to each other, no systematic connections were known until an insightful paper of Bernshteyn~\\cite{Bernshteyn2021LLL} who showed that results from distributed computing can automatically imply results in descriptive combinatorics.\nIn this work, we show that the connections between the three areas run much deeper than previously known, both in terms of techniques and in terms of complexity classes.\nIn fact, our work suggests that it is quite useful to consider all three perspectives as part of a common theory, and we will attempt to present our results accordingly.\nWe refer the reader to \\cref{fig:big_picture_trees} for a partial overview of the rich connections between the three perspectives, some of which are proven in this paper. \n\nIn this work, we focus on the case where the graph under consideration is a regular tree.\nDespite its simplistic appearance, regular trees play an important role in each of the three areas, as we will substantiate at the end of this section.\nTo already provide an example, in the area of distributed algorithms, almost all locality lower bounds are achieved on regular trees.\nMoreover, when regarding lower bounds, the property that they already apply on regular trees actually \\emph{strengthens} the result---a fact that is quite relevant for our work as our main contribution regarding the transfer of techniques between the areas is a new lower bound technique in the area of distributed computation that is an adaptation and generalization of a technique from descriptive combinatorics.\nRegarding our results about the relations between complexity classes from the three areas, we note that such connections are also studied in the context of paths and grids in two parallel papers \\cite{grebik_rozhon2021LCL_on_paths,grebik_rozhon2021toasts_and_tails}.\n\nIn the remainder of this section, we give a high-level overview of the three areas that we study.\nThe purpose of these overviews is to provide the reader with a comprehensive picture of the studied settings that can also serve as a starting point for delving deeper into selected topics in those areas---in order to follow our paper, it is not necessary to obtain a detailed understanding of the results and connections presented in the overviews. Necessary technical details will be provided in \\cref{sec:preliminaries}.\nMoreover, in \\cref{sec:contribution}, we present our contributions in detail.\n\n\n\n\n\n\\paragraph{Distributed Computing}\nThe definition of the $\\mathsf{LOCAL}$ model of distributed computing by Linial~\\cite{linial92LOCAL} was motivated by the desire to understand distributed algorithms in huge networks. \nAs an example, consider a huge network of wifi routers. Let us think of two routers as connected by an edge if they are close enough to exchange messages. \nIt is desirable that such close-by routers communicate with user devices on different channels to avoid interference.\nIn graph-theoretic language, we want to properly color the underlying network. \nEven if we are allowed a color palette with $\\Delta+1$ colors where $\\Delta$ denotes the maximum degree of the graph (which would admit a simple greedy algorithm in a sequential setting), the problem remains highly interesting in the distributed setting, as, ideally, each vertex decides on its output color after only a few rounds of communication with its neighbors, which does not allow for a simple greedy solution. \n\nThe $\\mathsf{LOCAL}$ model of distributed computing formalizes this setup: we have a large network, where each vertex knows the network's size, $n$, and perhaps some other parameters like the maximum degree $\\Delta$. In the case of randomized algorithms, each vertex has access to a private random bit string, while in the case of deterministic algorithms, each vertex is equipped with a unique identifier from a range polynomial in the size $n$ of the network. \nIn one round, each vertex can exchange any message with its neighbors and can perform an arbitrary computation. The goal is to find a solution to a given problem in as few communication rounds as possible. \nAs the allowed message size is unbounded, a $t$-round $\\mathsf{LOCAL}$ algorithm can be equivalently described as a function that maps $t$-hop neighborhoods to outputs---the output of a vertex is then simply the output its $t$-hop neighborhood\nmapped to by this function.\nAn algorithm is correct if and only if the collection of outputs at all vertices constitutes a correct solution to the problem. \n\nThere is a rich theory of distributed algorithms and the local complexity of many problems is understood. \nThe case of trees is a highlight of the theory: it is known that any local problem (a class of natural problems we will define later) belongs to one of only very few complexity classes.\nMore precisely, for any local problem, its randomized local complexity is either $O(1), \\Theta(\\log^* n), \\Theta(\\log\\log n), \\Theta(\\log n)$, or $\\Theta(n^{1\/k})$ for some $k \\in \\mathbb{N}$. Moreover, the deterministic complexity is always the same as the randomized one, except for the case $\\Theta(\\log\\log n)$, for which the corresponding deterministic complexity is $\\Theta(\\log n)$ (see \\cref{fig:big_picture_trees}). \n\n\n\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=\\textwidth]{fig\/big_picture_trees.pdf}\n \\caption{Complexity classes on regular trees considered in the three areas of distributed computing, factors of iid processes\/uniform algorithms, and descriptive combinatorics. The left part shows complexity classes of distributed computing. We use the shorthand $\\mathsf{LOCAL}$ if it does not matter whether we talk about the deterministic or randomized complexity. These two notions differ only for the class of problems of randomized local complexity $O(\\log\\log n)$, which have deterministic complexity $O(\\log n)$. \\\\\n The uniform complexity classes of sublogarithmic complexity are in correspondence to appropriate classes in the randomized local complexity model, as proven in \\cref{sec:separating_examples}. On the other hand, the class fiid is very similar to the class $\\mathsf{MEASURE}$ from descriptive combinatorics. The equivalence of the class $\\mathsf{CONTINUOUS}$ and $\\mathsf{LOCAL}(O(\\log^* n)) = \\mathsf{OLOCAL}(O(\\log^* 1\/\\varepsilon))$ is marked with a dashed arrow as it was proven in case the tree is generated by a group action (think of the tree being equipped with an additional $\\Delta$-edge coloring). The inclusion $\\mathsf{LOCAL}(O(\\log^* n)) \\subseteq \\mathsf{BOREL}$ however clearly holds also in our setting. The class $\\mathsf{BOREL}$ is incomparable with $\\mathsf{RLOCAL}(O(\\log\\log n))$, as proven in \\cref{sec:separating_examples}. }\n \\label{fig:big_picture_trees}\n\\end{figure}\n\n\n\\todo{thm 5.1 - 5.3}\n\n\n\n\n\\paragraph{(Finitary) Factors of iid Processes and Uniform Algorithms}\n\nIn recent years, \\emph{factors of iid (fiid) processes} on trees attracted a lot of attention in combinatorics, probability, ergodic theory and statistical physics \\cite{ABGGMP,ArankaFeriLaci,BackSzeg,BackVir,backhausz,backhausz2,BGHV,Ball,Bowen1,Bowen2,CPRR,csoka2015invariant,GabLyo,GamarnikSudan,GGP,HarangiVirag,HolroydPoisson,HLS,HPPS,HolroydPeres,HoppenWormald,Kun,RahmanPercolation,RahVir,LyonsTrees,Mester,schmidt2016circular,Timar1,Timar2}.\nIntuitively, factors of iid processes are randomized algorithms on, e.g., infinite $\\Delta$-regular trees, where each vertex outputs a solution to a problem after it explores random strings on vertices of the whole tree.\nAs an example, consider the \\emph{perfect matching} problem.\nAn easy parity argument shows that perfect matching cannot be solved by any local randomized algorithm on finite trees.\nHowever, if we allow a small fraction of vertices not to be matched, then, by a result of Nguyen and Onak \\cite{nguyen_onak2008matching} (see also \\cite{ElekLippner}), there is a constant-round randomized algorithm that produces such a matching on high-girth graphs (where the constant depends on the fraction of unmatched vertices that we allow).\nThis result can also be deduced from a result of Lyons and Nazarov \\cite{lyons2011perfect}, who showed that perfect matching can be described as a factor of iid process on an infinite $\\Delta$-regular tree.\nThe high-level idea behind this connection is that high-girth graphs approximate the infinite $\\Delta$-regular tree and constant-round local algorithms approximate factors of iid processes.\nThis correspondence is formalized in the notion of \\emph{Benjamini-Schramm} or \\emph{local-global} convergence \\cite{benjamini2011recurrence, hatamilovaszszegedy}.\nWe note that getting only ``approximate'' solutions, that is, solutions where a small fraction of vertices does not have to satisfy the constraints of a given problem, is intrinsic in this correspondence.\nRegardless, there are many techniques, such as entropy inequality \\cite{backhausz} or correlation decay \\cite{backhausz2}, and particular results such as the aforementioned perfect matching problem \\cite{lyons2011perfect} that provide lower and upper bounds, respectively, in our setting as well.\nWe refer the reader to \\cite{LyonsTrees,BGHV} for a comprehensive summary of the field. \n\n\nIn this paper, we mostly consider a stronger condition than fiid, namely so-called \\emph{finitary} factors of iid (ffiid) processes that are studied in the same context as fiid \\cite{holroyd2017one_dependent_coloring,HolroydSchrammWilson2017FinitaryColoring,Spinka}. \nPerhaps surprisingly, the notion of ffiid is identical to the notion of so-called uniform distributed randomized algorithms \\cite{Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring,grebik_rozhon2021toasts_and_tails} that we now describe. \nWe define an uniform local algorithm as a randomized local algorithm that does not know the size of the graph $n$ -- this enables us to run such an algorithm on infinite graphs, where there is no $n$. \nMore precisely, we require that each vertex eventually outputs a solution that is compatible with the output in its neighborhood, but the time until the vertex finishes is a potentially unbounded random variable. \nAs in the case of classical randomized algorithms, we can now measure the \\emph{uniform complexity} of an uniform local algorithm (known as the tail decay of ffiid \\cite{HolroydSchrammWilson2017FinitaryColoring}). \nThe uniform complexity of an algorithm is defined as the function $t(\\varepsilon)$ such that the probability that the algorithm run on a specific vertex needs to see outside its $t(\\varepsilon)$-hop neighborhood is at most $\\varepsilon$. \nAs in the case of classical local complexity, there is a whole hierarchy of possible uniform complexities (see \\cref{fig:big_picture_trees}). \n\nWe remark that uniform distributed local algorithms can be regarded as Las Vegas algorithms. The output will always be correct; there is however no fixed guarantee at what point all vertices have computed their final output. On the other hand, a randomized distributed local algorithm can be viewed as a Monte Carlo algorithm as it needs to produce an output after a fixed number of rounds, though the produced output might be incorrect. \n\n\n\n\n\n\n\n\n\\paragraph{Descriptive Combinatorics}\n \n \n\t\n\tThe Banach-Tarski paradox states that a three-dimensional ball of unit volume can be decomposed into finitely many pieces that can be moved by isometries (distance preserving transformations such as rotations and translations) to form two three-dimensional balls each of them with unit volume(!).\n The graph theoretic problem lurking behind this paradox is the following: fix finitely many isometries of $\\mathbb{R}^3$ and then consider a graph where $x$ and $y$ are connected if there is an isometry that sends $x$ to $y$.\n\tThen our task becomes to find a perfect matching in the appropriate subgraph of this graph -- namely, the bipartite subgraph where one partition contains points of the first ball and the other contains points of the other two balls. \n\tBanach and Tarski have shown that, with a suitably chosen set of isometries, the axiom of choice implies the existence of such a matching. \n\tIn contrast, since isometries preserve the Lebesgue measure, the pieces in the decomposition cannot be Lebesgue measurable. \n\tSurprisingly, Dougherty and Foreman \\cite{doughertyforeman} proved that the pieces in the Banach-Tarski paradox can have the \\emph{Baire property}. \n\tThe Baire property is a topological analogue of being Lebesgue measurable; a subset of $\\mathbb{R}^3$ is said to have the Baire property if its difference from some open set is topologically negligible.\n\t\n \n\t \n\t\n\t\n Recently, results similar to the Banach-Tarski paradox that lie on the border of combinatorics, logic, group theory, and ergodic theory led to an emergence of a new field often called \\textit{descriptive} or \\textit{measurable combinatorics}. \n The field focuses on the connection between the discrete and continuous and is largely concerned with the investigation of graph-theoretic concepts. \n The usual setup in descriptive combinatorics is that we have a graph with uncountably many connected components, each being a countable graph of bounded degree. \n For example, in case of the Banach-Tarski paradox, the vertices of the underlying graph are the points of the three balls, edges correspond to isometries, and the degree of each vertex is bounded by the number of chosen isometries. \n Some of the most beautiful results related to the field include \\cite{laczk, marksunger,measurablesquare,doughertyforeman,marks2016baire,gaboriau,KST,DetMarks,millerreducibility,conley2020borel,csokagrabowski,Bernshteyn2021LLL}, see \\cite{kechris_marks2016descriptive_comb_survey,pikhurko2021descriptive_comb_survey} for recent surveys. \n \n Importantly, in many of these results, including the Banach-Tarski paradox, graphs where each component is an infinite $\\Delta$-regular tree appear naturally. \n Oftentimes, questions considered in descriptive combinatorics lead to constructing a solution to a local problem in the underlying uncountable graph (in the case of Banach-Tarski, the local problem is perfect matching). \n The construction needs to be such that the solution of the problem has some additional regularity properties. For example in the case of Banach-Tarski, a solution is possible when the regularity condition is the Baire property, but not if it is Lebesgue measurability.\n In fact, together with Borel measurability these are the most prominent regularity conditions studied in descriptive combinatorics.\n \n The corresponding complexity classes of local problems that always admit a solution with the respective regularity property are $\\mathsf{BOREL},\\mathsf{MEASURE},\\mathsf{BAIRE}$ (See \\cref{fig:big_picture_trees}). \n In this paper, we moreover consider the setting where each connected component of the underlying graph is a $\\Delta$-regular tree. \n \n \n \n \n \n \n\t\n\t\n\t\t\n\t\n\t\t\n\t\t\n\n\t The connection between distributed computing and descriptive combinatorics arises from the fact that in descriptive combinatorics we care about constructions that do not use the axiom of choice. \n\t In the distributed language, the axiom of choice corresponds to leader election, that is, the constructions in descriptive combinatorics do not allow picking exactly one point in every component. \n\t \n\t To get some intuition about the power of the complexity class $\\mathsf{BOREL}$, we note that Borel constructions allow us to alternate countably many times the following two operations. \n\t First, any local algorithm with constant local complexity can be run. \n\t Second, we have an oracle that provides a maximal independent set (MIS) on any graph that can be constructed locally from the information computed so far \\cite{KST}. \n Note that from the speedup result of \\cite{chang_kopelowitz_pettie2019exp_separation} we get that every local problem with local complexity $O(\\log^* n)$ can be solved by constant local constructions and \\emph{one} call to such an MIS oracle. This implies the inclusion $\\mathsf{LOCAL}(O(\\log^* n)) \\subseteq \\mathsf{BOREL}$ in \\cref{fig:big_picture_trees} proven in the insightful paper of Bernshteyn \\cite{Bernshteyn2021LLL}. \n The relationship of the class $\\mathsf{MEASURE}$ (and of the class $\\mathsf{fiid}$ from the discussion of factors) to the class $\\mathsf{BOREL}$ is analogous to the relationship of randomized distributed algorithms to deterministic distributed algorithms. \n\n \n \n\n\n\n\n\n\n\t\n\n\t\n\t\n\t\n\n\n\n\n\n\n\\paragraph{Local Problems on Regular Trees}\n\nAfter introducing the three areas of interest in this work, we conclude the section by briefly discussing the kinds of problems we focus on, which are local problems on regular trees.\nMore precisely, we study \\emph{locally checkable labeling (LCL)} problems, which are a class of problems, where the correctness of the solution can be checked locally.\nExamples include classical problems from combinatorics such as proper vertex coloring, proper edge coloring, perfect matching, and maximal independent set.\nOne main goal of this paper is to understand possible complexity classes of LCLs without inputs on infinite $\\Delta$-regular trees and their finite analogues.\nWe refer the reader to \\cref{sec:preliminaries} for a precise definition of a finite $\\Delta$-regular tree.\n\nThe motivation for studying regular trees in this work stems from different sources: (a) infinite $\\Delta$-regular trees are studied in the area of ergodic theory \\cite{Bowen1,Bowen2}, random processes \\cite{backhausz,backhausz2,lyons2011perfect} and descriptive combinatorics \\cite{DetMarks,BrooksMeas}, (b) many lower bounds in distributed computing are proven in regular trees \\cite{balliu_et_al:binary_lcl,balliu2019LB,brandt_etal2016LLL,brandt19automatic,Brandt20tightLBmatching,ChangHLPU20,goeoes14_PODC}, and (c) connecting and comparing the techniques of the three areas in this simple setting reveals already deep connections, see \\cref{sec:contribution}.\n\n\n\n\n\n\n\n\n\n\\section{Our Contributions}\n\\label{sec:contribution}\n\nWe believe that our main contribution is presenting all three perspectives as part of a common theory. Our technical contribution is split into three main parts. \n\n\n\n\n\\subsection{Generalization of Marks' Technique}\\label{sec:intromarks}\nIn Section \\ref{sec:marks} we extend the Borel determinacy technique of Marks \\cite{DetMarks}, which was used to prove the nonexistence of Borel $\\Delta$-colorings and perfect matchings, to a broader class of problems, and \nadapt the extended technique to the distributed setting, thereby obtaining a simple method for proving distributed lower bounds.\nThis method is the first lower bound technique for distributed computing using ideas coming from descriptive combinatorics (see \\cite{BernshteynVizing} for a distributed computing upper bound motivated by descriptive combinatorics).\nMoreover, we show how to use the developed techniques to obtain both $\\mathsf{BOREL}$ and $\\mathsf{LOCAL}$ lower bounds for local problems from a natural class, called homomorphism problems.\nOur key technical ingredient for obtaining the mentioned techniques and results is a novel technique based on the notion of an \\emph{ID graph}.\nWe note that a very similar concept to the ID graph was independently discovered by \\cite{id_graph}. \n\n\\paragraph{Marks' technique}\nIn the following we give an introduction to Marks' technique by going through a variant of his proof~\\cite{DetMarks,Marks_Coloring} that shows that $\\Delta$-coloring has deterministic local complexity $\\Omega(\\log n)$. The proof already works in the case where the considered regular tree comes with an input $\\Delta$-edge coloring. In this case, the output color of a given vertex $u$ can be interpreted as that $u$ ``grabs'' the incident edge of that color. The problem is hence equivalent to the \\emph{edge grabbing} problem where every vertex is required to grab an incident edge such that no two vertices grab the same edge. \n\nWe first show the lower bound in the case that vertices do not have unique identifiers but instead are properly colored with $L > \\Delta$ colors. \nSuppose there is an algorithm $\\mathcal{A}$ solving the edge grabbing problem with local complexity $t(n) = o(\\log n)$, and consider a tree rooted at vertex $u$ of depth $t(n)$; such a tree has less than $n$ vertices, for large enough $n$.\nAssume that $u$ has input color $\\sigma \\in [L]$, and consider, for some fixed edge color $\\alpha$, the edge $e$ that is incident to $u$ and has color $\\alpha$.\nTwo players, Alice and Bob, are playing the following game. In the $i$-th round, Alice colors the vertices at distance $i$ from $u$ in the subtree reachable via edge $e$ with colors from $[L]$. Then, Bob colors all other vertices at distance $i$ from $u$ with colors from $[L]$ (see \\cref{fig:game}).\nConsider the output of $u$ when executing $\\mathcal{A}$ on the obtained colored tree.\nBob wins the game if $u$ grabs the edge $e$, and Alice wins otherwise. \n\nNote that either Alice or Bob has a winning strategy. \nGiven the color $\\sigma$ of $u$, if, for each edge color $\\alpha$, Alice has a winning strategy in the game corresponding to the pair $(\\sigma, \\alpha)$, then we can create $\\Delta$ copies of Alice and let them play their strategy on each subtree of $u$, telling them that the colors chosen by the other Alices are what Bob played. The result is a coloring of the input tree such that $u$, by definition, does not pick any edge, contradicting the fact that $\\mathcal{A}$ provides a valid solution!\nSo for every $\\sigma$ there is at least one $\\alpha$ such that Bob has a winning strategy for the game corresponding to $(\\sigma, \\alpha)$. By the pigeonhole principle, there are two colors $\\sigma_1, \\sigma_2$, such that Bob has a winning strategy for both pairs $(\\sigma_1, \\alpha)$ and $(\\sigma_2, \\alpha)$. \nBut now we can imagine a tree rooted in an edge between vertices $u_1, u_2$ that are colored with colors $\\sigma_1, \\sigma_2$. We can now take two copies of Bob, one playing at $u_1$ and the other playing at $u_2$ and let them battle it out, telling each copy that the other color from $\\{\\sigma_1, \\sigma_2\\}$ and whatever the other copy plays are the moves of Alice. The resulting coloring has the property that both $u_1$ and $u_2$, when executing $\\mathcal{A}$ on the obtained colored tree, grab the edge between them, a contradiction that finishes the proof!\n\n\\paragraph{The ID graph}\nThe downside of the proof is that it does not work in the model with unique identifiers (where the players' moves consist in assigning identifiers instead of colors), since gluing copies of the same player could result in an identifier assignment where the identifiers are not unique. \nOne possible remedy is to conduct the whole proof in the context of Borel graphs as was done by Marks. This proves an even stronger statement, namely that $\\Delta$-coloring is not in the class $\\mathsf{BOREL}$, but requires additional ad-hoc tricks and a pretty heavy set theoretic tool---Martin's celebrated Borel determinacy theorem~\\cite{martin} stating that even for infinite two-player games one of the players has to have a winning strategy if the payoff set is Borel. \nThe ID graph enables us to adapt the proof (and its generalization that we develop in \\cref{sec:marks}) to the distributed setting, where the fact that one of the players has a winning strategy is obvious. \nMoreover, we use an infinite version of the ID graph to generalize Marks' technique also in the Borel setting. \n\nHere is how it works: \nThe ID graph is a specific graph whose vertices are the possible unique input identifiers (for input graphs of size $n$), that is, numbers from $[n^{O(1)}]$. Its edges are colored with colors from $[\\Delta]$ and its girth is $\\Omega(\\log n)$. When we define the game between Alice and Bob, we require them to label vertices with identifiers in such a way that whenever a new vertex is labeled with identifier $i$, and its already labeled neighbor has identifier $j$, then $ij$ is an edge in the ID graph.\nMoreover, the color of edge $ij$ in the ID graph is required to be the same as the color of the edge between the vertices labeled $i$ and $j$ in the tree where the game is played. \nIt is straightforward to check that these conditions enforce that even if we let several copies of the same player play, the resulting tree is labeled with unique identifiers. Hence, the same argument as above now finally proves that the deterministic local complexity of $\\Delta$-coloring is $\\Omega (\\log n)$. \n\nWe note that our ID graph technique is of independent interest and may have applications in many different contexts. \nTo give an example from distributed computing, consider the proof of the deterministic $\\Omega(\\log n)$-round lower bound for $\\Delta$-coloring developed by the distributed community~\\cite{brandt_etal2016LLL,chang_kopelowitz_pettie2019exp_separation}, which is based on the celebrated round elimination technique. \nEven though the result is deterministic, the proof is quite technical due to the fact that it relies on examining randomized algorithms, for reasons similar to the reasons why Marks' proof does not apply directly to the setting with unique identifiers. \nFortunately, it can be again streamlined with the use of the ID graph technique. \nMoreover, in a follow-up work, the ID graph technique has already led to a new lower bound for the Lov\u00e1sz local lemma~\\cite{brandt_grunau_rozhon2021LLL_in_LCA} in the area of Local Computation Algorithms (which is part of the realm of sequential computation), thereby giving further evidence for the versatility of the technique. \n\n\n\\paragraph{Marks vs. Round Elimination}\nIt is quite insightful to compare Marks' technique (and our generalization of it) with the powerful round elimination technique~\\cite{brandt19automatic}, which has been responsible for all locality lower bounds of the last years~\\cite{brandt_etal2016LLL, balliu2019LB, brandt19automatic, balliu2019hardness_homogeneous, balliu_et_al:binary_lcl, Brandt20tightLBmatching, balliu20rulingset, ChangHLPU20, balliu21dominatingset}. \nWhile, on the surface, Marks' approach developed for the Borel world may seem quite different from the round elimination technique, there are actually striking similarities between the two methods.\nOn a high level, in the round elimination technique, the following argument is used to prove lower bounds in the LOCAL model:\nIf a $T$-round algorithm exists for a problem $\\Pi_0$ of interest, then there exists a $(T-1)$-round algorithm for some problem $\\Pi_1$ that can be obtained from $\\Pi_0$ in a mechanical manner.\nBy applying this step iteratively, we obtain a problem $\\Pi_t$ that can be solved in $0$ rounds; by showing that there is no $0$-algorithm for $\\Pi_t$ (which is easy to do if $\\Pi_t$ is known), a $(T+1)$-round lower bound for $\\Pi_0$ is obtained.\n\nThe interesting part regarding the relation to Marks' technique is how the ($T-i-1$)-round algorithms $\\mathcal{A}'$ are obtained from the ($T-i$)-round algorithms $\\mathcal{A}$ in the round elimination framework: in order to execute $\\mathcal{A}'$, each vertex $v$, being aware of its ($T-i-1$)-hop neighborhood, essentially asks whether, for all possible extensions of its view by one hop along a chosen incident edge, there exists some extension of its view by one hop along all other incident edges such that $\\mathcal{A}$, executed on the obtained ($T-i$)-hop neighborhood, returns a certain output at $v$, and then bases its output on the obtained answer.\nIt turns out that the vertex sets corresponding to these two extensions correspond precisely to two moves of the two players in the game(s) played in Marks' approach: more precisely, in round $T-i$ of a game corresponding to the considered vertex $v$ and the chosen incident edge, the move of Alice consists in labeling the vertices corresponding to the first extension, and the move of Bob consists in labeling the vertices corresponding to the second extension.\n\nHowever, despite the similarities, the two techniques (at least in their current forms) have their own strengths and weaknesses and are interestingly different in that there are local problems that we know how to obtain good lower bounds for with one technique but not the other, and vice versa.\nFinding provable connections between the two techniques is an exciting research direction that we expect to lead to a better understanding of the possibilities and limitations of both techniques.\n\nIn \\cref{sec:marks} we use our generalized and adapted version of Marks' technique to prove new lower bounds for so-called homomorphism problems. \nHomomorphism problems are a class of local problems that generalizes coloring problems---each vertex is to be colored with some color and there are constraints on which colors are allowed to be adjacent. \nThe constraints can be viewed as a graph---in the case of coloring this graph is a clique. \nIn general, whenever the underlying graph of the homomorphism problem is $\\Delta$-colorable, its deterministic local complexity is $\\Omega(\\log n)$, because solving the problem would imply that we can solve $\\Delta$-coloring too (in the same runtime). \nIt seems plausible that homomorphism problems of this kind are the only hard, i.e., $\\Omega(\\log n)$, homomorphism problems. However, our generalization of Marks' technique asserts that this is not true. \n\n\\begin{theorem}\nThere are homomorphism problems whose deterministic local complexity on trees is $\\Omega(\\log n)$ such that the chromatic number of the underlying graph is $2\\Delta-2$.\n\\end{theorem}\n\nIt is not known how to prove the same lower bounds using round elimination\\footnote{Indeed, the descriptions of the problems have comparably large numbers of labels and do not behave like so-called ``fixed points'' (i.e., nicely) under round elimination, which suggests that it is hard to find a round elimination proof with the currently known approaches.}; in fact, as far as we know, these problems are the only known examples of problems on $\\Delta$-regular trees for which a lower bound is known to hold but currently not achievable by round elimination. \nProving the same lower bounds via round elimination is an exciting open problem. \n\n\n\n\n\n\n\n\n\n\n\\subsection{Separation of Various Complexity Classes}\n\\label{sec:introseparation}\n\\paragraph{Uniform Complexity Landscape}\nWe investigate the connection between randomized and uniform distributed local algorithms, where uniform algorithms are equivalent to the studied notion of finitary factors of iid. \nFirst, it is simple to observe that local problems with uniform complexity $t(\\varepsilon)$ have randomized complexity $t(1 \/ n^{O(1)})$ -- by definition, every vertex knows its local output after that many rounds with probability $1 - 1\/n^{O(1)}$. The result thus follows by a union bound over the $n$ vertices of the input graph.\n\nOn the other hand, we observe that on $\\Delta$-regular trees the implication also goes in the opposite direction in the following sense. Every problem that has a randomized complexity of $t(n) = o(\\log n)$ has an uniform complexity of $O(t(1\/\\varepsilon))$. \n\nOne could naively assume that this equivalence also holds for higher complexities, but this is not the case. Consider for example the $3$-coloring problem.\nIt is well-known in the distributed community that $3$-coloring a tree can be solved deterministically in $O(\\log n)$ rounds using the rake-and-compress decomposition~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup,MillerR89}.\nOn the other hand, there is no uniform algorithm for $3$-coloring a tree. \nIf there were such an uniform algorithm, we could run it on any graph with large enough girth and color $99\\%$ of its vertices with three colors. This in turn would imply that the high-girth graph has an independent set of size at least $0.99 \\cdot n\/3$. This is a contradiction with the fact that there exists high-girth graphs with a much smaller independence number~\\cite{bollobas}. \n\nInterestingly, the characterization of Bernshteyn \\cite{Bernshteyn_work_in_progress} implies that \\emph{any} uniform distributed algorithm can be ``sped up'' to a deterministic local $O(\\log n)$ complexity, as we prove in~\\cref{thm:MainBaireLog}. \n\nWe show that there are local problems that can be solved by an uniform algorithm but only with a complexity of $\\Omega(\\log 1\/\\varepsilon)$. Namely, the problem of constructing a $2$-hop perfect matching on infinite $\\Delta$-regular trees for $\\Delta \\ge 3$ has an uniform local complexity between $\\Omega(\\log 1\/\\varepsilon)$ and $O(\\operatorname{poly}\\log 1\/\\varepsilon)$. \nFormally, this proves the following theorem. \n\n\\begin{theorem}\n$\\mathsf{OLOCAL}(O(\\log\\log 1\/\\varepsilon)) \\subsetneq \\mathsf{OLOCAL}(O(\\operatorname{poly}\\log 1\/\\varepsilon))$. \n\\end{theorem}\n\nThe uniform algorithm for this problem is based on a so-called one-ended forest decomposition introduced in \\cite{BrooksMeas} in the descriptive combinatorics context. In a one-ended forest decomposition, each vertex selects exactly one of its neighbors as its parent by orienting the corresponding edge outwards. This defines a decomposition of the vertices into infinite trees. We refer to such a decomposition as a one-ended forest decomposition if the subtree rooted at each vertex only contains finitely many vertices. \nHaving computed such a decomposition, $2$-hop perfect matching can be solved inductively starting from the leaf vertices of each tree.\n\nWe leave the understanding of the uniform complexity landscape in the regime $\\Omega(\\log 1\/\\varepsilon)$ as an exciting open problem. In particular, does there exist a function $g(\\varepsilon)$ such that each local problem that can be solved by an uniform algorithm has an uniform complexity of $O(g(\\varepsilon))$?\n\n\n\n\\paragraph{Relationship of Distributed Classes with Descriptive Combinatorics}\n\nBernshteyn recently proved that $\\mathsf{LOCAL}(\\log^* n) \\subseteq \\mathsf{BOREL}$ \\cite{Bernshteyn2021LLL}. That is, each local problem with a deterministic $\\mathsf{LOCAL}$ complexity of $O(\\log^* n)$ also admits a Borel-measurable solution. A natural question to ask is whether the converse also holds. \nIndeed, it is known that $\\mathsf{LOCAL}(\\log^* n) = \\mathsf{BOREL}$ on paths with no additional input \\cite{grebik_rozhon2021LCL_on_paths}. \nWe show that on regular trees the situation is different. \nOn one hand, a characterization of Bernshteyn \\cite{Bernshteyn_work_in_progress} implies that $\\mathsf{BOREL} \\subseteq \\mathsf{BAIRE} \\subseteq \\mathsf{LOCAL}(O(\\log n))$. \nOn the other hand, we show that this result cannot be strengthened by proving the following result. \n\n\\begin{theorem}\n$\\mathsf{BOREL} \\subsetneq \\mathsf{RLOCAL}(o(\\log n))$.\n\\end{theorem}\n\nThat is, there exists a local problem that admits a Borel-measurable solution but cannot be solved with a (randomized) $\\mathsf{LOCAL}$ algorithm running in a sublogarithmic number of rounds. \n\nLet us sketch a weaker separation, namely that $\\mathsf{BOREL} \\setminus \\mathsf{LOCAL}(O(\\log^* n)) \\neq \\emptyset$. Consider a version of $\\Delta$-coloring where a subset of vertices can be left uncolored. However, the subgraph induced by the uncolored vertices needs to be a collection of doubly-infinite paths (in finite trees, this means each path needs to end in a leaf vertex). \nThe nonexistence of a fast distributed algorithm for this problem essentially follows from the celebrated $\\Omega(\\log n)$ deterministic lower bound for $\\Delta$-coloring of~\\cite{brandt_etal2016LLL}. \nOn the other hand, the problem allows a Borel solution. First, sequentially compute $\\Delta-2$ maximal independent sets, each time coloring all vertices in the MIS with the same color, followed by removing all the\ncolored vertices from the graph. In that way, a total of $\\Delta-2$ colors are used. Moreover, each\nuncolored vertex has at most $2$ uncolored neighbors. \nThis implies that the set of uncolored vertices forms a disjoint union of finite paths, one ended\ninfinite paths and doubly infinite paths. The first two classes can be colored inductively with two\nadditional colors, starting at one endpoint of each path in a Borel way (namely it can be done by making use of the countably many MISes in larger and larger powers of the input graph). \nHence, in the end only doubly infinite paths are left uncolored, as desired. \n\nTo show the stronger separation between the classes $\\mathsf{BOREL}$ and $\\mathsf{RLOCAL}(o(\\log n))$ we use a variation of the\n$2$-hop perfect matching problem. In this variation, some of the vertices can be left unmatched, but similar as in the variant of the $\\Delta$-coloring problem described above, the graph induced by all the unmatched vertices needs to satisfy some additional constraints.\n\nWe conclude the paragraph by noting that the separation between the classes $\\mathsf{BOREL}$ and $\\mathsf{LOCAL}(O(\\log^* n))$ is not as simple as it may look in the following sense. \nThis is because problems typically studied in the $\\mathsf{LOCAL}$ model with a $\\mathsf{LOCAL}$ complexity of $\\omega(\\log^* n)$ like $\\Delta$-coloring and perfect matching also do not admit a Borel-measurable solution due to the technique of Marks \\cite{DetMarks} that we discussed in \\cref{sec:intromarks}.\n\n\n\n\n\n\\subsection{$\\mathsf{LOCAL}(O(\\log n)) = \\mathsf{BAIRE}$} \n\nWe already discussed that one of complexity classes studied in descriptive combinatorics is the class $\\mathsf{BAIRE}$.\nRecently, Bernshteyn proved \\cite{Bernshteyn_work_in_progress} that all local problems that are in the complexity class $\\mathsf{BAIRE}, \\mathsf{MEASURE}$ or $\\mathsf{fiid}$ have to satisfy a simple combinatorial condition which we call being $\\ell$-full. \nOn the other hand, all $\\ell$-full problems allow a $\\mathsf{BAIRE}$ solution \\cite{Bernshteyn_work_in_progress}.\nThis implies a complete combinatorial characterization of the class $\\mathsf{BAIRE}$. \nWe defer the formal definition of $\\ell$-fullness to Section \\ref{sec:baire} as it requires a formal definition of a local problem. \nInformally speaking, in the context of vertex labeling problems, a problem is $\\ell$-full if we can choose a subset $S$ of the labels with the following property.\nWhenever we label two endpoints of a path of at least $\\ell$ vertices with two labels from $S$, we can extend the labeling with labels from $S$ to the whole path such that the overall labeling is valid.\nFor example, proper $3$-coloring is $3$-full with $S=\\{1,2,3\\}$ because for any path of three vertices such that its both endpoints are colored arbitrarily, we can color the middle vertex so that the overall coloring is proper. \nOn the other hand, proper $2$-coloring is not $\\ell$-full for any $\\ell$. \n\nWe complement this result as follows. \nFirst, we prove that any $\\ell$-full problem has local complexity $O(\\log n)$, thus proving that all complexity classes considered in the areas of factors of iids and descriptive combinatorics from \\cref{fig:big_picture_trees} are contained in $\\mathsf{LOCAL}(O(\\log n))$. \nIn particular, this implies that the existence of \\emph{any} uniform algorithm implies a local distributed algorithm for the same problem of local complexity $O(\\log n)$. \nWe obtain this result via the well-known rake-and-compress decomposition~\\cite{MillerR89}. \n\n\nOn the other hand, we prove that any problem in the class $\\mathsf{LOCAL}(O(\\log n))$ satisfies the $\\ell$-full condition. \nThe proof combines a machinery developed by Chang and Pettie~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}\nwith additional nontrivial ideas.\nIn this proof we construct recursively a sequence of sets of rooted, layered, and partially labeled trees, where the partial labeling is computed by simulating any given $O(\\log n)$-round distributed algorithm, \nand then the set $S$ meeting the $\\ell$-full condition is constructed by considering all possible extensions of the partial labeling to complete correct labeling of these trees.\n\nThis result implies the following equality:\n\\begin{theorem}\n$\\mathsf{LOCAL}(O(\\log n)) = \\mathsf{BAIRE}$. \n\\end{theorem}\n\nThis equality is surprising in that the definitions of the two classes do not seem to have much in common on the first glance!\nMoreover, the proof of the equality relies on nontrivial results in both distributed algorithms (the technique of Chang and Pettie~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}) and descriptive combinatorics (the fact that a hierarchical decomposition, so-called toast, can be constructed in $\\mathsf{BAIRE}$, \\cite{conleymillerbound}, see Proposition \\ref{pr:BaireToast}). \n\n\n\n\n\n\n\n\n\n\n\n\nThe combinatorial characterization of the local complexity class $\\mathsf{LOCAL}(O(\\log n))$ on $\\Delta$-regular trees is interesting from the perspective of distributed computing alone. \nThis result can be seen as a part of a large research program aiming at classification of possible local complexities on various graph classes~\\cite{balliu2020almost_global_problems,brandt_etal2016LLL,chang_kopelowitz_pettie2019exp_separation,chang_pettie2019time_hierarchy_trees_rand_speedup,chang2020n1k_speedups,Chang_paths_cycles,balliu2021_rooted_trees,balliu_et_al:binary_lcl,balliu2018new_classes-loglog*-log*}. \nThat is, we wish not only to understand possible complexity classes (see the left part of \\cref{fig:big_picture_trees} for possible local complexity classes on regular trees), but also to find combinatorial characterizations of problems in those classes that allow us to efficiently decide for a given problem which class it belongs to. \nUnfortunately, even for grids with input labels, it is \\emph{undecidable} whether a given local problem can be solved in $O(1)$ rounds~\\cite{naorstockmeyer,brandt_grids}, since\nlocal problems on grids can be used to simulate a Turing machine.\nThis undecidability result does not apply to paths and trees, hence for these graph classes it is still hopeful that we can find simple and useful characterizations for different classes of distributed problems. \n\nIn particular, on paths it is decidable what classes a given local problem belongs to, for all classes coming from the three areas considered here, and this holds even if we allow inputs~\\cite{Chang_paths_cycles,grebik_rozhon2021LCL_on_paths}. \nThe situation becomes much more complicated when we consider trees. \nRecently, classification results on trees were obtained for so-called binary-labeling problems \\cite{balliu_et_al:binary_lcl}. \nMore recently, a complete classification was obtained in the case of \\emph{rooted regular trees}~\\cite{balliu2021_rooted_trees}. \nAlthough their algorithm takes exponential time in the worst case, the authors provided a practical implementation fast enough to classify many typical problems of interest. \n\nMuch less is known for general, \\emph{unoriented} trees, with an arbitrary number of labels. \nIn general, deciding the optimal distributed complexity for a local problem on bounded-degree trees is $\\mathsf{EXPTIME}$-hard~\\cite{chang2020n1k_speedups}, such a hardness result does not rule out the possibility for having a simple and polynomial-time characterization for the case of \\emph{regular trees}, where there is no input and the constraints are placed only on degree-$\\Delta$ vertices. \nIndeed, it was stated in~\\cite{balliu2021_rooted_trees} as an open question to find such a characterization. Our characterization of $\\mathsf{LOCAL}(O(\\log n)) = \\mathsf{BAIRE}$ by $\\ell$-full problems makes progress in better understanding the distributed complexity classes on trees and towards answering this open question. \n\n\n\n\n\n\n\n\\paragraph{Roadmap}\n\\hypertarget{par:roadmap}{}\nIn \\cref{sec:preliminaries}, we define formally all the three setups we consider in the paper. \nIn \\cref{sec:marks} we discuss the lower bound technique of Marks and the new concept of an ID graph.\nNext, in \\cref{sec:separating_examples} we prove some basic results about the uniform complexity classes and give examples of problems separating some classes from \\cref{fig:big_picture_trees}. \nFinally, in \\cref{sec:baire} we prove that a problem admits a Baire measurable solution if and only if it admits a distributed algorithm of local complexity $O(\\log n)$.\n\nThe individual sections can be read largely independently of each other.\nMoreover, most of our results that are specific to only one of the three areas can be understood without reading the parts of the paper that concern the other areas.\nWe encourage the reader interested mainly in one of the areas to skip the respective parts.\n\n\n\n\n\\section{Preliminaries}\n\\label{sec:preliminaries}\n\nIn this section, we explain the setup we work with, the main definitions and results. \nThe class of graphs that we consider in this work are either infinite $\\Delta$-regular trees, or their finite analogue that we define formally in \\cref{subsec:local_problems}.\n\nWe sometimes explicitly assume $\\Delta > 2$. \nThe case $\\Delta = 2$, that is, studying paths, behaves differently and seems much easier to understand \\cite{grebik_rozhon2021LCL_on_paths}. \nUnless stated otherwise, we do not consider any additional structure on the graphs, but sometimes it is natural to work with trees with an input $\\Delta$-edge-coloring. \n\n\n\\subsection{Local Problems on $\\Delta$-regular trees}\n\\label{subsec:local_problems}\n\nThe problems we study in this work are locally checkable labeling (LCL) problems, which, roughly speaking, are problems that can be described via local constraints that have to be satisfied in a suitable neighborhood of each vertex.\nIn the context of distributed algorithms, these problems were introduced in the seminal work by Naor and Stockmeyer~\\cite{naorstockmeyer}, and have been studied extensively since.\nIn the modern formulation introduced in~\\cite{brandt19automatic}, instead of labeling vertices or edges, LCL problems are described by labeling half-edges, i.e., pairs of a vertex and an incident edge.\nThis formulation is very general in that it not only captures vertex and edge labeling problems, but also others such as orientation problems, or combinations of all of these types.\nBefore we can provide this general definition of an LCL, we need to introduce some definitions.\nWe start by formalizing the notion of a half-edge.\n\n\\begin{definition}[Half-edge]\n\\label{def:half-edge}\n A \\emph{half-edge} is a pair $(v, e)$ where $v$ is a vertex, and $e$ an edge incident to $v$.\n We say that a half-edge $(v, e)$ is \\emph{incident} to a vertex $w$ if $w = v$, we say that a vertex $w$ is \\emph{contained} in a half edge $(v,e)$ if $w=v$, and we say that $(v, e)$ \\emph{belongs} to an edge $e'$ if $e' = e$.\n We denote the set of all half-edges of a graph $G$ by $H(G)$.\n A \\emph{half-edge labeling} is a function $c \\colon H(G) \\to \\Sigma$ that assigns to each half-edge an element from some label set $\\Sigma$.\n\\end{definition}\n\nIn order to speak about finite $\\Delta$-regular trees, we need to consider slightly modified definition of a graph.\nWe think of each vertex to be contained in $\\Delta$-many half-edges, however, not every half edge belongs to an actual edge of the graph.\nThat is half-edges are pairs $(v,e)$, but $e$ is formally not a pair of vertices.\nSometimes we refer to these half-edges as \\emph{virtual} half-edges.\nWe include a formal definition to avoid confusions. See also \\cref{fig:Delta_reg_tree}. \n\n\n\\begin{definition}[$\\Delta$-regular trees]\n A tree $T$, finite or infinite is a \\emph{$\\Delta$-regular tree} if either it is infinite and $T=T_\\Delta$, where $T_\\Delta$ is the unique infinite $\\Delta$-regular tree, that is each vertex has exactly $\\Delta$-many neighbors, or it is finite of maximum degree $\\Delta$ and each vertex $v\\in T$ of degree $d\\le \\Delta$ is contained in $(\\Delta-d)$-many virtual half-edges.\n \n \n Formally, we can view $T$ as a triplet $(V(T), E(T), H(T))$, where $(V(T),E(T))$ is a tree of maximum degree $\\Delta$ and $H(T)$ consists of real half-edges, that is pairs $(v,e)$, where $v\\in V(T)$, $e\\in E(T)$ and $e$ is incident to $v$, together with some virtual edges, in the case when $T$ is finite, such that each vertex $v\\in V(T)$ is contained in exactly $\\Delta$-many half-edges (real or virtual).\n \n \n\\end{definition}\n\n\\begin{figure}\n \\centering\n \\includegraphics{fig\/Delta_reg_tree.pdf}\n \\caption{A $3$-regular tree on 15 vertices. Every vertex has $3$ half-edges but not all half-edges lead to a different vertex. }\n \\label{fig:Delta_reg_tree}\n\\end{figure}\n\nAs we are considering trees in this work, each LCL problem can be described in a specific form that provides two lists, one describing all label combinations that are allowed on the half-edges incident to a vertex, and the other describing all label combinations that are allowed on the two half-edges belonging to an edge.\\footnote{Every problem that can be described in the form given by Naor and Stockmeyer~\\cite{naorstockmeyer} can be equivalently described as an LCL problem in this list form, by simply requiring each output label on some half-edge $h$ to encode all output labels in a suitably large (constant) neighborhood of $h$ in the form given in~\\cite{naorstockmeyer}.}\nWe arrive at the following definition for LCLs on $\\Delta$-regular trees.\\footnote{Note that the defined LCL problems do not allow the correctness of the output to depend on input labels.}\n\n\\begin{definition}[LCLs on $\\Delta$-regular trees]\\label{def:LCL_trees}\nA \\emph{locally checkable labeling problem}, or \\emph{LCL} for short, is a triple $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$, where $\\Sigma$ is a finite set of labels, $\\mathcal{V}$ is a subset of unordered cardinality-$\\Delta$ multisets\\footnote{Recall that a multiset is a modification of the concept of sets, where repetition is allowed.} of labels from $\\Sigma$ , and $\\mathcal{E}$ is a subset of unordered cardinality-$2$ multisets of labels from $\\Sigma$.\n\nWe call $\\mathcal{V}$ and $\\mathcal{E}$ the \\emph{vertex constraint} and \\emph{edge constraint} of $\\Pi$, respectively.\nMoreover, we call each multiset contained in $\\mathcal{V}$ a \\emph{vertex configuration} of $\\Pi$, and each multiset contained in $\\mathcal{E}$ an \\emph{edge configuration} of $\\Pi$.\n\nLet $T$ be a $\\Delta$-regular tree and $c:H(T)\\to \\Sigma$ a half-edge labeling of $T$ with labels from $\\Sigma$.\nWe say that $c$ is a \\emph{$\\Pi$-coloring}, or, equivalently, a correct solution for $\\Pi$, if, for each vertex $v$ of $T$, the multiset of labels assigned to the half-edges incident to $v$ is contained in $\\mathcal{V}$, and, for each edge $e$ of $T$, the cardinality-$2$ multiset of labels assigned to the half-edges belonging to $e$ is an element of $\\mathcal{E}$.\n\\end{definition}\n\nAn equivalent way to define our setting would be to consider $\\Delta$-regular trees as commonly defined, that is, there are vertices of degree $\\Delta$ and vertices of degree $1$, i.e., leaves.\nIn the corresponding definition of LCL one would consider leaves as unconstrained w.r.t. the vertex constraint, i.e., in the above definition of a correct solution the condition ``for each vertex $v$'' is replaced by ``for each non-leaf vertex $v$''.\nEquivalently, we could allow arbitrary trees of maximum degree $\\Delta$ as input graphs, but, for vertices of degree $< \\Delta$, we require the multiset of labels assigned to the half-edges to be extendable to some cardinality-$\\Delta$ multiset in $\\mathcal{V}$.\nWhen it helps the exposition of our ideas and is clear from the context, we may make use of these different but equivalent perspectives.\n\nWe illustrate the difference between our setting and ``standard setting'' without virtual half-edges on the perfect matching problem.\nA standard definition of the perfect matching problem is that some edges are picked in such a way that each vertex is covered by exactly one edge.\nIt is easy to see that there is no local algorithm to solve this problem on the class of finite trees (without virtual half-edges), this is a simple parity argument.\nHowever, in our setting, every vertex needs to pick exactly one half-edge (real or virtual) in such a way that both endpoints of each edge are either picked or not picked.\nWe remark that in our setting it is not difficult to see that (if $\\Delta>2$), then this problem can be solved by a local deterministic algorithm of local complexity $O(\\log(n))$.\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{The $\\mathsf{LOCAL}$ model}\n\nIn this section, we define local algorithms and local complexity.\nWe discuss the general relation between the classical local complexity and the uniform local complexity. Recall that when we talk about distributed algorithm on $\\Delta$-regular trees, the algorithm has access to $n$, the size of the tree. The measure of complexity is the classical local complexity. On the other hand, when we talk about uniform distributed algorithms on $\\Delta$-regular trees, we talk about an infinite $\\Delta$-regular tree and the measure of complexity is the uniform local complexity. \nWe start with the classical notions.\nRecall that $B(v,t)$ is the $t$-hop neighborhood of a vertex $v$ in an underlying graph $G$.\n\n\\begin{definition}[Local algorithm]\n\\label{def:local_algorithm}\nA \\emph{distributed local algorithm} $\\mathcal{A}$ of local complexity $t(n)$ is a function defined on all possible $t(n)$-hop neighborhoods of a vertex. Applying an algorithm $\\mathcal{A}$ on an input graph $G$ means that the function is applied to a $t(n)$-hop neighborhood $B(u, t(n))$ of each vertex $u$ of $G$. The output of the function is a labeling of the half-edges around the given vertex. The algorithm also takes as input the size of the input graph $n$. \n\\end{definition}\n\n\\begin{definition}[Local complexity]\n\\label{def:local_complexity}\nWe say that an LCL problem $\\Pi$ has a deterministic local complexity $t(n)$ if there is a local algorithm $\\mathcal{A}$ of local complexity $t(n)$ such that when run on the input graph $G$, with each of its vertices having a unique identifier from $[n^{O(1)}]$, $\\mathcal{A}$ always returns a valid solution to $\\Pi$. \nWe also say $\\Pi \\in \\mathsf{DLOCAL}( O(t(n))$. \n\nThe randomized complexity is define analogously, but instead of unique identifiers, each vertex of $G$ has an infinite random string. The solution of $\\mathcal{A}$ needs to be a valid solution with probability $1 - 1\/n$. \nWe also say $\\Pi \\in \\mathsf{RLOCAL}( O(t(n))$. \n\\end{definition}\n\nWhenever we talk about \\emph{local} complexity on $\\Delta$-regular trees, we always tacitly think about the class of finite $\\Delta$-regular trees. We use the notation $\\mathsf{LOCAL}(O(t(n)))$ whenever the deterministic and the randomized complexity of the problems are the same up to constant factor (cf. \\cref{thm:basicLOCAL}). \n\n\\paragraph{Classification of Local Problems on $\\Delta$-regular Trees}\n\nThere is a lot of work aiming to classify all possible local complexities of LCL problems on bounded degree trees. This classification was recently finished.\nEven though the results in the literature are stated for the classical notion of finite trees, we note that it is easy to check that the same picture emerges if we restrict ourselves to $\\Delta$-regular trees according to our definition.\n\n\\begin{theorem}[Classification of local complexities of LCLs on $\\Delta$-regular trees \\cite{naorstockmeyer,chang_kopelowitz_pettie2019exp_separation,chang_pettie2019time_hierarchy_trees_rand_speedup,chang2020n1k_speedups,balliu2020almost_global_problems,balliu2019hardness_homogeneous,brandt_grunau_rozhon2021classification_of_lcls_trees_and_grids}]\n\\label{thm:basicLOCAL}\nLet $\\Pi$ be an LCL problem. Then the deterministic\/randomized local complexity of $\\Pi$ on $\\Delta$-regular trees is one of the following:\n\\begin{enumerate}\n \\item $O(1)$,\n \\item $\\Theta(\\log^*n)$,\n \\item $\\Theta(\\log n)$ deterministic and $\\Theta(\\log\\log n)$ randomized,\n \\item $\\Theta(\\log n)$,\n \\item $\\Theta(n^{1\/k})$ for $k \\in \\mathbb{N}$.\n\\end{enumerate}\n\\end{theorem}\n\nWe do not know whether the complexity classes $\\Theta(n^{1\/k})$ are non-empty on \\emph{$\\Delta$-regular} trees. The $2 \\frac{1}{2}$ coloring problems $\\mathcal{P}_k$ are known to have complexity $\\Theta(n^{1\/k})$ on \\emph{bounded-degree} trees~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}. However, as the correctness criterion of $\\mathcal{P}_k$ depends on the degree of vertices, it does not fit into the class of LCL problems that we consider in \\cref{def:LCL_trees}. \nNevertheless, as can be seen in \\cref{fig:big_picture_trees}, all classes of factors of iid and from descriptive combinatorics are already contained in the class $\\mathsf{LOCAL}(O(\\log n))$. \n\n\n\n\\paragraph{Uniform Algorithms}\nAs we mentioned before, when talking about local complexities, we always have in mind that the underlying graph is finite.\nIn particular, the corresponding algorithm knows the size of the graph.\nOn infinite $\\Delta$-regular trees, or infinite graphs in general, we use the following notion \\cite{HolroydSchrammWilson2017FinitaryColoring,Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring}. \n\n\\begin{definition}\nAn \\emph{uniform local algorithm} $\\mathcal{A}$ is a function that is defined on all possible (finite) neighborhoods of a vertex.\nFor some neighborhoods it outputs a special symbol $\\emptyset$ instead of a labeling of the half-edges around the central vertex.\nApplying $\\mathcal{A}$ on a graph $G$ means that for each vertex $u$ of $G$ the function is applied to $B(u,t)$, where $t$ is the minimal number such that $\\mathcal{A}(u,t)\\not=\\emptyset$.\nWe call $t$ the coding radius of $\\mathcal{A}$, and denote it, as a function on vertices, as $R_\\mathcal{A}$.\n\\end{definition}\n\nWe define the corresponding notion of \\emph{uniform local complexity} for infinite $\\Delta$-regular trees where each vertex is assigned an infinite random string.\n\n\\begin{definition}[Uniform local complexity \\cite{HolroydSchrammWilson2017FinitaryColoring}]\n\\label{def:uniform_complexity}\nWe say that the \\emph{uniform local (randomized) complexity} of an LCL problem $\\Pi$ is $t(\\varepsilon)$ if there is an uniform local algorithm $\\mathcal{A}$ such that the following hold on the infinite $\\Delta$-regular tree.\nRecall that $R_\\mathcal{A}$ is the random variable measuring the coding radius of a vertex $u$, that is, the distance $\\mathcal{A}$ needs to look at to decide the answer for $u$. Then, for any $0 < \\varepsilon < 1$:\n\\[\n\\P(R_\\mathcal{A} \\ge t(\\varepsilon)) \\leq \\varepsilon.\n\\]\nWe also say $\\Pi \\in \\mathsf{OLOCAL}(O(t(\\varepsilon))$.\n\\end{definition}\n\nWe finish by stating the following lemma that bounds the uniform complexity of concatenation of two uniform algorithm (we need to be little bit more careful and cannot just add the complexites up). \n\n\\begin{lemma}[Sequential Composition]\n\\label{lem:sequential_composition}\nLet $A_1$ and $A_2$ be two distributed uniform algorithms with an uniform local complexity of $t_1(\\varepsilon)$ and $t_2(\\varepsilon)$, respectively. \nLet $A$ be the sequential composition of $A_1$ and $A_2$. That is, $A$ needs to know the output that $A_2$ computes when the local input at every vertex is equal to the output of the algorithm $A_1$ at that vertex. Then, $t_A(\\varepsilon) \\leq t_{A_1} \\left( \\frac{\\varepsilon\/2}{\\Delta^{t_{A_2}(\\varepsilon\/2) + 1}}\\right) + t_{A_2}(\\varepsilon\/2)$.\n\\end{lemma}\n\\begin{proof}\nConsider some arbitrary vertex $u$. Let $E_1$ denote the event that the coding radius of $\\mathcal{A}_1$ is at most $t_{A_1} \\left( \\frac{\\varepsilon\/2}{\\Delta^{t_{A_2}(\\varepsilon\/2) + 1}}\\right)$ for all vertices in the $t_{A_2}(\\varepsilon\/2)$-hop neighborhood around $u$. As the $t_{A_2}(\\varepsilon\/2)$-hop neighborhood around $u$ contains at most $\\Delta^{t_{A_2}(\\varepsilon\/2) + 1}$ vertices, a union bound implies that $P(E_1) \\geq 1 - \\varepsilon\/2$. Moreover, let $E_2$ denote the event that the coding radius of algorithm $\\mathcal{A}_2$ at vertex $u$ is at most $t_{A_2}(\\varepsilon\/2)$. By definition, $P(E_2) \\geq 1 - \\varepsilon\/2$. Moreover, if both events $E_1$ and $E_2$ occur, which happens with probability at least $1 - \\varepsilon$, then the coding radius of algorithm $\\mathcal{A}$ is at most $t_{A_1} \\left( \\frac{\\varepsilon\/2}{\\Delta^{t_{A_2}(\\varepsilon\/2) + 1}}\\right) + t_{A_2}(\\varepsilon\/2)$, thus finishing the proof.\n \\end{proof}\n\n\n\\subsection{Descriptive combinatorics}\n\\label{subsec:descriptivecombinatorics}\n\nBefore we define formally the descriptive combinatorics complexity classes, we give a high-level overview on their connection to distributing computing for the readers more familiar with the latter.\n\nThe complexity class that partially captures deterministic local complexity classes is called $\\mathsf{BOREL}$ (see also Remark \\ref{rem:cont}).\nFirst note that by a result of Kechris, Solecki and Todor\\v{c}evi\\'c \\cite{KST} the \\emph{maximal independent set problem} (with any parameter $r\\in \\mathbb{N}$) is in this class for any bounded degree graph.\\footnote{That is, it is possible to find a Borel maximal independent set, i.e., a maximal independent set which is, moreover, a Borel subset of the vertex set.}\nIn particular, this yields that $\\mathsf{BOREL}$ contains the class $\\mathsf{LOCAL}(O(\\log^* n))$ by the characterization of \\cite{chang_kopelowitz_pettie2019exp_separation}, see \\cite{Bernshteyn2021LLL}. Moreover, as mentioned before, $\\mathsf{BOREL}$ is closed under countably many iterations of the operations of finding maximal independent set (for some parameter that might grow) and of applying a constant local rule that takes into account what has been constructed previously.\\footnote{It is in fact an open problem, whether this captures fully the class $\\mathsf{BOREL}$. However, note that an affirmative answer to this question would yield that problems can be solved in an ``effective\" manner in the Borel context, which is known not to be the case in unbounded degree graphs \\cite{todorvcevic2021complexity}.}\n\nTo get a better grasp of what this means, consider for example the proper vertex $2$-coloring problem on half-lines. It is clear that no local algorithm can solve this problem.\nHowever, as it is possible to determine the starting vertex after countably many iterations of the maximal independent set operation, we conclude that this problem is in the class $\\mathsf{BOREL}$.\nThe idea that $\\mathsf{BOREL}$ can compute some unbounded, global, information will be implicitly used in all the construction in \\cref{sec:separating_examples} that separate $\\mathsf{BOREL}$ from local classes.\n\nThe intuition behind the class $\\mathsf{MEASURE}$ is that it relates in the same way to the class $\\mathsf{BOREL}$, as randomized local algorithms relate to deterministic ones.\nIn particular, the operations that are allowed in the class $\\mathsf{MEASURE}$ are the same as in the class $\\mathsf{BOREL}$ but the solution of a given LCL can be incorrect on a measure zero set.\n\nThe class $\\mathsf{BAIRE}$ can be considered as a topological equivalent of the measure theoretic class $\\mathsf{MEASURE}$, that is, a solution can be incorrect on a topologically negligible set.\nThe main difference between the classes $\\mathsf{MEASURE}$ and $\\mathsf{BAIRE}$ is that in the later there is a hierarchical decomposition that is called \\emph{toast}.\n(Note that this phenomenon is present in the case of $\\mathsf{MEASURE}$ exactly on so-called amenable graphs. It is also tightly connected with the notion of hyperfiniteness \\cite{conleymillerbound,gao2015forcing}.)\nThe independence of colorings on a tree together with this structure allows for a combinatorial characterization of the class $\\mathsf{BAIRE}$, which was proven by Bernshteyn \\cite{Bernshteyn_work_in_progress}, see also \\cref{sec:baire}.\n\n\n\n\nNext we formulate the precise definitions.\nWe refer the reader to \\cite{pikhurko2021descriptive_comb_survey,kechris_marks2016descriptive_comb_survey,Bernshteyn2021LLL,kechrisclassical}, or to \\cite[Introduction,~Section~4.1]{grebik_rozhon2021LCL_on_paths} and \\cite[Section~7.1,~7.2]{grebik_rozhon2021toasts_and_tails} for intuitive examples and standard facts of descriptive set theory.\nIn particular, we do not define here the notion standard Borel\/probability space, a Polish topology, a Borel probability measure, Baire property etc.\n\nLet $\\mathcal{G}$ be a Borel graph of bounded maximum degree on a standard Borel space $X$.\nIn this paper we consider exclusively acyclic $\\Delta$-regular Borel graphs and we refer to them as \\emph{$\\Delta$-regular Borel forests}.\nIt is easy to see that the set of half-edges (see Definition \\ref{def:half-edge}) is naturally a standard Borel space, we denote this set by $H(\\mathcal{G})$.\nThus, it makes sense to consider Borel labelings of $H(\\mathcal{G})$.\nMoreover, if $\\mathcal{G}$ is a $\\Delta$-regular Borel forest and $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ is an LCL, we can also decide whether a coloring $f:\\mathcal{H}(\\mathcal{G})\\to \\Sigma$ is a solution to $\\Pi$ as in Definition \\ref{def:LCL_trees}. \nSimilarly, we say that the coloring $f$ solves $\\Pi$, e.g., on a $\\mu$-conull set for some Borel probability measure $\\mu$ on $X$ if there is a Borel set $C\\subseteq X$ such that $\\mu(C)=1$, the vertex constraints are satisfied around every $x\\in C$ and the edge constraints are satisfied for every $x,y\\in C$ that form an edge in $\\mathcal{G}$.\n\n\\begin{definition}[Descriptive classes]\n\\label{def:descriptive}\nLet $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ be an LCL.\nWe say that $\\Pi$ is in the class $\\mathsf{BOREL}$ if for every acyclic $\\Delta$-regular Borel graph $\\mathcal{G}$ on a standard Borel space $X$, there is a Borel function $f:H(\\mathcal{G})\\to \\Sigma$ that is a $\\Pi$-coloring of $\\mathcal{G}$.\n\nWe say that $\\Pi$ is in the class $\\mathsf{BAIRE}$ if for every acyclic $\\Delta$-regular Borel graph $\\mathcal{G}$ on a standard Borel space $X$ and every compatible Polish topology $\\tau$ on $X$, there is a Borel function $f:H(\\mathcal{G})\\to \\Sigma$ that is a $\\Pi$-coloring of $\\mathcal{G}$ on a $\\tau$-comeager set.\n\nWe say that $\\Pi$ is in the class $\\mathsf{MEASURE}$ if for every acyclic $\\Delta$-regular Borel graph $\\mathcal{G}$ on a standard Borel space $X$ and every Borel probability measure $\\mu$ on $X$, there is a Borel function $f:H(\\mathcal{G})\\to \\Sigma$ that is a $\\Pi$-coloring of $\\mathcal{G}$ on a $\\mu$-conull set.\n\\end{definition}\n\nThe following claim follows directly from the definition.\n\n\\begin{claim}\\label{cl:basicDC}\nWe have $\\mathsf{BOREL}\\subseteq \\mathsf{MEASURE}, \\mathsf{BAIRE}$.\n\\end{claim}\n\nRecently Bernshteyn \\cite{Bernshteyn2021LLL,Bernshteyn2021local=cont} proved several results that connect distributed computing with descriptive combinatorics.\nUsing the language of complexity classes we can formulate some of the result as inclusions in \\cref{fig:big_picture_trees}.\n\n\\begin{theorem}[\\cite{Bernshteyn2021LLL}]\nWe have\n\\begin{itemize}\n \\item $\\mathsf{LOCAL}(O(\\log^*(n)))\\subseteq \\mathsf{BOREL}$,\n \\item $\\mathsf{RLOCAL}(o(\\log(n))) \\subseteq \\mathsf{MEASURE}, \\mathsf{BAIRE}$.\n\\end{itemize}\nIn fact, these inclusions hold for any reasonable class of bounded degree graphs.\n\\end{theorem}\n\n\\begin{remark}\n\\label{rem:cont}\nThe ``truly'' local class in descriptive combinatorics is the class $\\mathsf{CONTINUOUS}$.\nEven though we do not define this class here\\footnote{To define the class $\\mathsf{CONTINUOUS}$, rather than asking for a continuous solution on all possible $\\Delta$-regular Borel graphs, one has to restrict to a smaller family of graphs, otherwise the class trivializes. To define precisely this family is somewhat inconvenient, and not necessary for our arguments.}, and we refer the reader to \\cite{Bernshteyn2021local=cont,GJKS,grebik_rozhon2021toasts_and_tails} for the definition and discussions in various cases, we mention that the inclusion \\ \\[\\mathsf{LOCAL}(O(\\log^* n))\\subseteq \\mathsf{CONTINUOUS} \\tag{*}\\] holds in most reasonable classes of bounded degree graphs, see \\cite{Bernshteyn2021LLL}.\nThis also applies to our situation\nRecently, it was shown by Bernshteyn \\cite{Bernshteyn2021local=cont}, and independently by Seward \\cite{Seward_personal}, that $(*)$ can be reversed for Cayley graphs of finitely generated groups.\nThis includes, e.g., natural $\\Delta$-regular trees with proper edge $\\Delta$-coloring, as this is a Cayely graph of free product of $\\Delta$-many $\\mathbb{Z}_2$ induced by the standard generating set.\nIt is, however, not clear whether it $(*)$ can be reversed in our situation, i.e., $\\Delta$-regular trees without any additional labels.\n\n\n\n\\end{remark}\n\n\n\n\n\n\n\\subsection{Random processes}\n\nWe start with an intutitve description of fiid processes. Let $T_\\Delta$ be the infinite $\\Delta$-regular tree.\nInformally, \\emph{factors of iid processes (fiid)} on $T_\\Delta$ are infinite analogues of local randomized algorithms in the following way.\nLet $\\Pi$ be an LCL and $u\\in T_\\Delta$.\nIn order to solve $\\Pi$, we are allowed to explore the whole graph, and the random strings that are assigned to vertices, and then output a labeling of half-edges around $u$.\nIf such an assignment is a measurable function and produces a $\\Pi$-coloring almost surely, then we say that $\\Pi$ is in the class $\\mathsf{fiid}$.\nMoreover, if every vertex needs to explore only finite neighborhood to output a solution, then we say that $\\Pi$ is in the class $\\mathsf{ffiid}$.\nSuch processes are called \\emph{finitary} fiid (\\emph{ffiid}). There is also an intimate connection between ffiid and uniform algorithms.\nThis is explained in \\cite[Section 2.2]{grebik_rozhon2021toasts_and_tails}.\nInformally, an ffiid process that almost surely solves $\\Pi$ is, in the language of distributed computing, an uniform local algorithm that solves $\\Pi$.\nThis allows us to talk about uniform local complexity of an ffiid.\nIn the rest of the paper we interchange both notions freely with slight preference for the distributed computing language.\n\nNow we define formally these classes, using the language of probability.\nWe denote by $\\operatorname{Aut}(T_\\Delta)$ the automorphism group of $T_\\Delta$.\nAn \\emph{iid process} on $T_\\Delta$ is a collection of iid random variables $Y=\\{Y_v\\}_{v\\in V(T_\\Delta)}$ indexed by vertices, or edges, half-edges etc, of $T_\\Delta$ such that their joint distribution is invariant under $\\operatorname{Aut}(T_\\Delta)$.\nWe say that $X$ is a \\emph{factor of iid process (fiid)} if $X=F(Y)$, where $F$ is a measurable $\\operatorname{Aut}(T_\\Delta)$-invariant map and $Y$ is an iid process on $T_\\Delta$.\\footnote{We always assume $Y=(2^\\mathbb{N})^{T_\\Delta}$ endowed with the product Lebesgue measure.}\nMoreover, we say that $X$ is a \\emph{finitary factor if iid process (ffiid)} if $F$ depends with probability $1$ on a finite (but random) radius around each vertex.\nWe denote as $R_F$ the random variable that assigns minimal such radius to a given vertex, and call it the \\emph{coding radius} of $F$.\nWe denote the space of all $\\Pi$-colorings of $T_\\Delta$ as $X_\\Pi$.\nThis is a subspace of $\\Sigma^{H(T_\\Delta)}$ that is invariant under $\\operatorname{Aut}(T_\\Delta)$.\n\n\\begin{definition}\nWe say that an LCL $\\Pi$ is in the class $\\mathsf{fiid}$ ($\\mathsf{ffiid}$) if there is an fiid (ffiid) process $X$ that almost surely produces elements of $X_\\Pi$.\n\nEquivalently, we can define $\\mathsf{ffiid} = \\bigcup_f \\mathsf{OLOCAL}(f(\\varepsilon))$ where $f$ ranges over all functions. That is, $\\mathsf{ffiid}$ is the class of problems solvable by any uniform distributed algorithm. \n\n\\end{definition}\n\n\nIt is obvious that $\\mathsf{ffiid}\\subseteq \\mathsf{fiid}$.\nThe natural connection between descriptive combinatorics and random processes is formulated by the inclusion $\\mathsf{MEASURE}\\subseteq \\mathsf{fiid}$.\nWhile this inclusion is trivially satisfied, e.g., in the case of $\\Delta$-regular trees with proper edge $\\Delta$-coloring, in our situation we need a routine argument that we include for the sake of completeness in \\cref{app:ktulu}.\n\n\\begin{restatable}{lemma}{FIIDinMEASURE}\\label{thm:FIIDinMEASURE}\nLet $\\Pi$ be an LCL such that $\\Pi\\in \\mathsf{MEASURE}$.\nThen $\\Pi\\in \\mathsf{fiid}$.\n\\end{restatable}\n\n\n\n\n\\subsection{Specific Local Problems}\n\nHere we list some well-known local problems that were studied in all three considered areas. We describe the results known about these problems with respect to classes from \\cref{fig:big_picture_trees}. \n\n\\paragraph{Edge Grabbing}\nWe start by recalling the well-known problem of edge grabbing (a close variant of the problem is known as sinkless orientation\\cite{brandt_etal2016LLL}). \nIn this problem, every vertex should mark one of its half-edges (that is, grab an edge) in such a way that no edge can be marked by two vertices. \nIt is known that $\\Pi_{\\text{edgegrab}} \\not\\in \\mathsf{LOCAL}(O(\\log^* n))$ \\cite{brandt_etal2016LLL} but $\\Pi_{\\text{edgegrab}} \\in \\mathsf{RLOCAL}(O(\\log\\log n))$.\n\nSimilarly, $\\Pi_{\\text{edgegrab}} \\not \\in \\mathsf{BOREL}$ by \\cite{DetMarks}, but $\\Pi_{\\text{edgegrab}}\\in \\mathsf{MEASURE}$ by \\cite{BrooksMeas}: to see the former, just note that if a $\\Delta$-regular tree admits proper $\\Delta$-colorings of both edges and vertices, every vertex can grab an edge of the same color as the color of the vertex. Thus, $\\Pi_{\\text{edgegrab}}\\in \\mathsf{BOREL}$ would yield that $\\Delta$-regular Borel forests with Borel proper edge-colorings admit a Borel proper $\\Delta$-coloring, contradicting \\cite{DetMarks}.\n\nThis completes the complexity characterization of $\\Pi_{\\text{edgegrab}}$ as well as the proper vertex $\\Delta$-coloring with respect to classes in \\cref{fig:big_picture_trees}. \n\n\n\n\n\n\\paragraph{Perfect Matching}\nAnother notorious LCL problem, whose position in \\cref{fig:big_picture_trees} is, however, not completely understood, is the perfect matching problem $\\Pi_{\\text{pm}}$.\nRecall that the perfect matching problem $\\Pi_{\\text{pm}}$ asks for a matching that covers all vertices of the input tree.\\footnote{In our formalism this means that around each vertex exactly one half-edge is picked.}\nIt is known that $\\Pi_{\\text{pm}} \\in \\mathsf{fiid}$ \\cite{lyons2011perfect}, and it is easy to see that $\\Pi_{\\text{pm}} \\not\\in \\mathsf{RLOCAL}(O(\\log\\log(n)))$ (we will prove a stronger result in \\cref{pr:AddingLineLLL}).\nMarks proved \\cite{DetMarks} that it is not in $\\mathsf{BOREL}$, even when the underlying tree admits a Borel bipartition.\nIt is not clear if $\\Pi_{\\text{pm}}$ is in $\\mathsf{ffiid}$, nor whether it is in $\\mathsf{MEASURE}$.\n\n\n\n\n\n\\paragraph{Graph Homomorphism}\nWe end our discussion with LCLs that correspond to graph homomorphisms (see also the discussion in Section \\ref{sec:marks}).\nThese are also known as edge constraint LCLs.\nLet $G$ be a finite graph.\nThen we define $\\Pi_G$ to be the LCL that asks for a homomorphism from the input tree to $G$, that is, vertices are labeled with vertices of $G$ in such a way that edge relations are preserved.\nThere are not many positive results except for the class $\\mathsf{BAIRE}$.\nIt follows from the result of Bernshteyn \\cite{Bernshteyn_work_in_progress} (see \\cref{sec:baire}) that $\\Pi_G\\in \\mathsf{BAIRE}$ if and only if $G$ is not bipartite.\nThe main negative results can be summarized as follows.\nAn easy consequence of the result of Marks \\cite{DetMarks} is that if $\\chi(G)\\le \\Delta$, then $\\Pi_G\\not \\in \\mathsf{BOREL}$.\nIn this paper, we describe examples of graphs of chromatic number up to $2\\Delta-2$ such that the corresponding homomorphism problem is not in $\\mathsf{BOREL}$, see \\cref{sec:marks}.\nThe theory of \\emph{entropy inequalities} see \\cite{backhausz} implies that if $G$ is a cycle on more than $9$ vertices, then $\\Pi_G\\not\\in \\mathsf{fiid}$.\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Generalization of Marks' technique}\n\\label{sec:marks}\n\nIn this section, we first develop a new way of proving lower bounds in the $\\mathsf{LOCAL}$ model based on a generalization of a technique of Marks \\cite{DetMarks}. Then, we use ideas arising in the finitary setting---connected to the adaptation of Marks' technique---to obtain new results back in the Borel context.\nFor an introduction to Marks' technique and a high-level discussion about the challenges in adapting the technique to the standard distributed setting as well as our solution via the new notion of an ID graph, we refer the reader to \\cref{sec:intromarks}.\n\n\\begin{comment}\nIn \\cite{DetMarks}, Marks showed that there exist $\\Delta$-regular Borel forests that do not admit a Borel proper $\\Delta$-coloring. As it was indicated above, this is a striking difference between $\\mathsf{MEASURE}$ or $\\mathsf{BAIRE}$ and $\\mathsf{BOREL}$, as such colorings exist in the former classes, see \\cite{BrooksMeas}. Since prior to Marks' breakthrough (essentially) any lower bound on Borel chromatic numbers utilized measure theoretic or Baire category arguments, his approach was (and had to be) based on a new technique: \\emph{determinacy of infinite games}. \n\nFix a payoff set $A \\subseteq \\mathbb{R}$. Consider the infinite game in which two players alternatingly determine digits of a real number. The first player wins iff the resulting real is in $A$. A rather counter-intuitive fact of set theory is that there exists a set $A$ for which none of the players has a winning strategy, i.e., the game is not determined. However, a deep theorem of Martin \\cite{martin} says that if $A$ is Borel, then one of the players does have such a strategy. \nMarks has constructed $\\Delta$-regular Borel forests and games with Borel payoff sets, for which the existence of a winning strategy for either of the players contradicts the existence of a Borel proper $\\Delta$-coloring. \n\nIn what follows, we generalize and adapt his games to the finitary setting, thus obtaining new lower bounds for LOCAL algorithms; then we use ideas inspired by the LOCAL technique of ID graphs to prove a new theorem in the Borel context. \n\nNote that Marks considered Borel versions of certain locally checkable problems (e.g., coloring, edge grabbing, perfect matching). As discussed above, for these problems, lower bounds were independently obtained via the round elimination technique in the LOCAL context, see e.g., \\cite{brandt19automatic}\\todo{cite more stuff??-Z}. Remarkably, the core ideas behind Brandt's and Marks' argument exhibit a surprising degree of similarity. \n\n\\todo{add the story about going from borel to local, and then back from local to borel}\n\nHowever, our extension of Marks' technique yields lower bounds for fairly natural problems that were resistant to round elimination.\nMoreover, as with our approach the lower bound is obtained by checking whether the considered LCL satisfies a certain combinatorial property called \\emph{playability condition} (see Definition~\\ref{def:Playability}), we can use it to obtain lower bounds for problems that have complicated descriptions---something that the round elimination technique in its current form is ill-suited for. \\todo{Maybe this sentence should rather be somewhere in the intro?}\n\n\n\n\\end{comment}\n\n\n\n\n\n\n\nThe setting we consider in this section is $\\Delta$-regular trees that come with a proper $\\Delta$-edge coloring with colors from $[\\Delta]$.\nAll lower bounds discussed already hold under these restrictions (and therefore also in any more general setting).\n\nRecall that an LCL $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ is given by specifying a list of allowed vertex configurations and a list of allowed edge configurations (see Definition \\ref{def:LCL_trees}).\nTo make our lower bounds as widely applicable as possible, we replace the latter list by a separate list for each of the $\\Delta$ edge colors; in other words, we consider LCLs where the correctness of the output is allowed to depend on the input that is provided by the edge colors.\nHence, formally, in this section, an LCL is more generally defined: it is a tuple $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ where $\\Sigma$ and $\\mathcal{V}$ are as before, while $\\mathcal{E} = (\\mathcal{E}_{\\alpha})_{\\alpha \\in [\\Delta]}$ is now a $\\Delta$-tuple of sets $\\mathcal{E}_{\\alpha}$ consisting of cardinality-$2$ multisets.\nSimilarly as before, a half-edge labeling (see Definition \\ref{def:descriptive} for the Borel context) with labels from $\\Sigma$ is a correct solution for $\\Pi$ if, for each vertex $v$, the multiset of labels assigned to the half-edges incident to $v$ is contained in $\\mathcal{V}$, and, for each edge $e$, the multiset of labels assigned to the half-edges belonging to $e$ is contained in $\\mathcal{E}_{\\alpha}$, where $\\alpha$ is the color of the edge $e$.\n\nThe idea of our approach is to identify a condition that an LCL $\\Pi$ necessarily has to satisfy if it is solvable in $O(\\log^* n)$ rounds in the $\\mathsf{LOCAL}$ model.\nShowing that $\\Pi$ does not satisfy this condition then immediately implies an $\\Omega(\\log n)$ deterministic and $\\Omega(\\log \\log n)$ randomized lower bound, by Theorem~\\ref{thm:basicLOCAL}.\n\nIn order to define our condition we need to introduce the notion of a configuration graph: a \\emph{configuration graph} is a $\\Delta$-tuple $\\mathbb{P}=(\\mathbb{P}_\\alpha)_{\\alpha\\in \\Delta}$ of graphs, where the vertex set of each of the graphs $\\mathbb{P}_\\alpha$ is the set of subsets of $\\Sigma$, and there is an edge in $\\mathbb{P}_\\alpha$ connecting two vertices $S,T$ if and only if there are $\\mathtt{a}\\in S$ and $\\mathtt{b}\\in T$ such that $\\{\\mathtt{a},\\mathtt{b}\\}\\in {\\mathcal{E}}_\\alpha$, note that loops are allowed.\n(Naturally, we will consider any two vertices of different $\\mathbb{P}_\\alpha$ to be distinct, even if they correspond to the same subset of $\\Sigma$.)\n\nNow we are set to define the aforementioned condition. The intuition behind the playability condition is the following: assume that there exists a local algorithm $\\mathcal{A}$ that solves $\\Pi$ using the $t$ neighbourhood of a given vertex. We are going to define a family of two player games. The game will depend on some $S \\subseteq \\Sigma$. $\\operatorname{Alice}$ and $\\operatorname{Bob}$ will assign labels (or IDs) to the vertices in the $t$-neighbourhood (in some way specified later on, depending on $\\alpha \\in [\\Delta]$). When the assignment is complete, we evaluate $\\mathcal{A}$ on the root of the obtained labelled graph, this way obtaining an element $s$ of $\\Sigma$, and decide who is the winner based on $s \n\\in S$ or not. Naturally, it depends on $S$ and $\\alpha$, which player has a winning strategy. This gives rise to a two coloring of vertices of $\\mathbb{P}_\\alpha$ by colors $\\operatorname{Alice}$ and $\\operatorname{Bob}$. The failure of the playability condition will guarantee that using a strategy stealing argument one can derive a contradiction.\n\n\n\n\n\\begin{definition}[Playability condition]\\label{def:Playability}\nWe say that an LCL $\\Pi$ is \\emph{playable} if for every $\\alpha\\in [\\Delta]$ there is a coloring $\\Lambda_{\\alpha}$ of the vertices of $\\mathbb{P}_{\\alpha}$ with two colors $\\{\\operatorname{Alice},\\operatorname{Bob}\\}$ such that the following conditions are satisfied:\n\\begin{enumerate}\n \\item [(A)] For any tuple $(S_{\\alpha})_{\\alpha \\in [\\Delta]} \\in V(\\mathbb{P}_1) \\times \\dots \\times V(\\mathbb{P}_{\\Delta})$ satisfying $\\Lambda_\\alpha(S_\\alpha)=\\operatorname{Alice}$ for each $\\alpha \\in [\\Delta]$, there exist $\\mathtt{a}_\\alpha\\not\\in S_\\alpha$ such that $\\{\\mathtt{a}_\\alpha\\}_{\\alpha\\in \\Delta}\\in {\\mathcal{V}}$, and\n \\item [(B)] for any $\\alpha \\in [\\Delta]$, and any tuple $(S,T) \\in V(\\mathbb{P}_{\\alpha}) \\times V(\\mathbb{P}_{\\alpha})$ satisfying $\\Lambda_\\alpha(S)=\\Lambda_\\alpha(T)=\\operatorname{Bob}$, we have that $(S,T)$ is an edge of $\\mathbb{P}_{\\alpha}$.\n\\end{enumerate}\n\n\\end{definition}\n\nOur aim in this section is to show the following general results.\n\n\\begin{theorem}\\label{th:MainLOCAL}\nLet $\\Pi$ be an LCL that is not playable.\nThen $\\Pi$ is not in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\n\\end{theorem}\n\nUsing ideas from the proof of this result, we can formulate an analogous theorem in the Borel context.\nLet us mention that while \\cref{th:MainLOCAL} is a consequence of \\cref{thm:MainBorel} by \\cite{Bernshteyn2021LLL}, we prefer to state the theorems separately.\nThis is because the proof of the Borel version of the theorem uses heavily complex results such as the Borel determinacy theorem, theory of local-global limits and some set-theoretical considerations about measurable sets.\nThis is in a stark contrast with the proof of \\cref{th:MainLOCAL} that uses `merely' the existence of large girth graphs and determinacy of finite games.\nHowever, the ideas surrounding these concepts are the same in both cases.\n\n\\begin{theorem}\\label{thm:MainBorel}\nLet $\\Pi$ be an LCL that is not playable\nThen $\\Pi$ is not in the class $\\mathsf{BOREL}$.\n\\end{theorem}\n\nThe main application of these abstract results is to find lower bounds for \\emph{(graph) homomorphism LCLs} aka \\emph{edge constraint LCLs}.\nWe already defined this class in \\cref{subsec:local_problems} but we recall the definition in the formalism of this section to avoid any confusion.\nLet $G$ be a finite graph.\nThen $\\Pi_G=(\\Sigma,\\mathcal{V},\\mathcal{E})$ is defined by letting\n\\begin{enumerate}\n\\item \\label{p:vertex} ${\\mathcal{V}}=\\{({\\mathtt{a}},\\dots,{\\mathtt{a}}):\\mathtt{a}\\in \\Sigma\\}$,\n\\item $\\Sigma=V(G)$,\n\\item $\\forall \\alpha \\in [\\Delta] \\ (\\mathcal{E}_\\alpha=E(G)).$\n\\end{enumerate}\nAn LCL for which \\eqref{p:vertex} holds is called a \\emph{vertex-LCL}. There is a one-to-one correspondence between vertex-LCLs for which $\\forall \n\\alpha,\\beta \\ (\\mathcal{E}_\\alpha=\\mathcal{E}_\\beta)$ and LCLs of the form $\\Pi_G$. (Indeed, to vertex-LCLs with this property one can associate a graph $G$ whose vertices are the labels in $\\Sigma$ and where two vertices $\\ell, \\ell'$ are connected by an edge if and only if $\\{ \\ell, \\ell' \\} \\in \\mathcal{E}$.)\nNote that if $\\Pi$ is a vertex-LCL, then condition (A) in Definition~\\ref{def:Playability} is equivalent to the statement that no tuple $\\{S_\\alpha\\}_{\\alpha\\in \\Delta}$ that satisfies the assumption of (A) can cover $\\Sigma$, i.e., that $\\Sigma \\nsubseteq S_1 \\cup \\dots \\cup S_{\\Delta}$ for such a tuple.\nMoreover, if $\\Pi_G$ is a homomorphism LCL, then $\\mathbb{P}_\\alpha=\\mathbb{P}_{\\beta}$ for every $\\alpha,\\beta\\in \\Delta$.\n\nThe simplest example of a graph homomorphism LCL is given by setting $G$ to be the clique on $k$ vertices; the associated LCL $\\Pi_G$ is simply the problem of finding a proper $k$-coloring of the vertices of the input graph.\nNow we can easily derive Marks' original result from our general theorem.\n\n\\begin{corollary}[Marks \\cite{DetMarks}]\\label{cor:detmarks}\nLet $G$ be a finite graph that has chromatic number at most $\\Delta$. Then $\\Pi_G$ is not playable. In particular, there is a $\\Delta$-regular Borel forest that have Borel chromatic number $\\Delta+1$. \n\\end{corollary}\n\\begin{proof}\nLet $A_1,\\dots, A_\\Delta$ be some independent sets that cover $G$, and $\\Lambda_1, \\dots, \\Lambda_{\\Delta}$ arbitrary colorings of the vertices of $\\mathbb{P}_1, \\dots, \\mathbb{P}_{\\Delta}$, respectively, with colors from $\\{\\operatorname{Alice},\\operatorname{Bob}\\}$.\nIt follows that $\\Lambda_\\alpha(A_\\alpha)=\\operatorname{Alice}$ for every $\\alpha\\in \\Delta$, since otherwise condition (B) in Definition~\\ref{def:Playability} is violated with $S=T=A_\\alpha$.\nBut then condition (A) does not hold.\n\\end{proof}\n\nAs our main application we describe graphs with chromatic number larger than $\\Delta$ such that $\\Pi_G$ is not playable.\nThis rules out the hope that the complexity of $\\Pi_G$ is connected with chromatic number being larger than $\\Delta$.\nIn Section \\ref{subsec:homlcls} we show the following.\nNote that the case $k=\\Delta$ is \\cref{cor:detmarks}.\n\n\\begin{theorem}\n\\label{th:homomorphism}\nLet $\\Delta > 2$ and $\\Delta < k \\leq 2\\Delta-2$. There exists a graph $G_k$ with $\\chi(G_k)= k$, such that $\\Pi_{G_k}$ is not playable and $\\Pi_{G_k} \\in \\mathsf{RLOCAL}(O(\\log\\log n ))$. In particular, $\\Pi_{G_k}\\not\\in \\mathsf{BOREL}$ and $\\Pi_{G_k}\\not\\in\\mathsf{LOCAL}(O(\\log^*(n)))$.\n\\end{theorem}\n\nInterestingly, recent results connected to counterexamples to Hedetniemi's conjecture yield the same statement asymptotically, as $\\Delta \\to \\infty$ (see Remark \\ref{r:hedet}).\n\n\\begin{remark}\nIt can be shown that for $\\Delta=3$ both the Chv\\' atal and Gr\\\" otsch graphs are suitable examples for $k=4$.\n\\end{remark}\n\n\\begin{remark}\nAs another application of his technique Marks showed in \\cite{DetMarks} that there is a Borel $\\Delta$-regular forest that does not admit Borel perfect matching.\nThis holds even when we assume that the forest is Borel bipartite, i.e., it has Borel chromatic number $2$.\nIn order to show this result Marks works with free joins of colored hyperedges, that is, Cayley graphs of free products of cyclic groups.\nOne should think of two types of triangles ($3$-hyperedges) that are joined in such a way that every vertex is contained in both types and there are no cycles.\nWe remark that the playability condition can be formulated in this setting.\nSimilarly, one can derive a version of \\cref{thm:MainBorel}.\nHowever, we do not have any application of this generalization.\n\\end{remark}\n\n\n\n\n\n\n\n\n\\subsection{Applications of playability to homomorphism LCLs}\\label{subsec:homlcls}\n\nIn this section we find the graph examples from \\cref{th:homomorphism}.\nFirst we introduce a condition $\\Delta$-(*) that is a weaker condition than having chromatic number at most $\\Delta$, but still implies that the homomorphism LCL is not playable. Then, we will show that--similarly to the way the complete graph on $\\Delta$-many vertices, $K_\\Delta$, is maximal among graphs of chromatic number $\\leq \\Delta$--there exists a maximal graph (under homorphisms) with property $\\Delta$-(*). Recall that we assume $\\Delta >2$.\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=.6\\textwidth]{fig\/h3.pdf}\n \\caption{The maximal graph with the property $\\Delta$-(*) for $\\Delta=3$.\n \\label{fig:h3}\n\\end{figure}\n\n\n\n\n\\begin{definition}[Property $\\Delta$-(*)]\nLet $\\Delta>2$ and $G=(V,E)$ be a finite graph.\nWe say that $G$ satisfies property $\\Delta$-(*) if there are sets $S_0,S_1\\subseteq V$ such that $G$ restricted to $V \\setminus S_i$ has chromatic number at most $(\\Delta-1)$ for $i\\in \\{0,1\\}$, and there is no edge between $S_0$ and $S_1$.\n\\end{definition}\n\nNote that $\\chi(G) \\leq \\Delta$ implies $\\Delta$-(*): indeed, if $A_1,\\dots,A_\\Delta$ are independent sets that cover $V(G)$, we can set $S_0=S_1=A_1$.\n\nOn the other hand, we claim that if $G$ satisfies $\\Delta$-(*) then $\\chi(G)\\le 2\\Delta-2$. In order to see this, take $S_0,S_1 \\subseteq V(G)$ witnessing $\\Delta$-(*). Then, as there is no edge between $S_0$ and $S_1$, so in particular, between $S_0 \\setminus S_1$ and $S_0 \\cap S_1$, it follows that the chromatic number of $G$'s restriction to $S_0$ is $\\leq \\Delta-1$. But then we can construct proper $\\Delta-1$-colorings of $S_0$ and $V(G) \\setminus S_0$, which shows our claim.\n\n\n\\begin{proposition}\\label{pr:CONdition*}\nLet $G$ be a graph satisfying $\\Delta$-(*).\nThen $\\Pi_G$ is not playable.\n\\end{proposition}\n\\begin{proof}\nFix $S_0,S_1$ as in the definition of $\\Delta$-(*) and assume for a contradiction that colorings $\\Lambda_1, \\dots, \\Lambda_\\Delta$ as described in Definition~\\ref{def:Playability} exist.\nBy Property $\\Delta$-(*), there exist independent sets $A_1, \\dots, A_{\\Delta-1}$ such that $S_0$ together with the $A_i$ covers $G$, i.e., such that $S_0 \\cup A_1 \\cup \\dots \\cup A_{\\Delta-1} = V(G)$.\nFor each $\\alpha \\in [\\Delta-1]$, we must have $\\Lambda_\\alpha(A_{\\alpha})=\\operatorname{Alice}$, since otherwise condition (B) in Definition~\\ref{def:Playability} is violated.\nConsequently $\\Lambda_{\\Delta}(S_0)=\\operatorname{Bob}$, otherwise condition (A) is violated.\nSimilarly $\\Lambda_\\Delta(S_1)=\\operatorname{Bob}$.\nThis shows that $\\Lambda_\\Delta$ does not satisfy condition (B) with $S = S_0$ and $T = S_1$.\n\\end{proof}\n\nNext we describe maximal examples of graphs that satisfy the condition $\\Delta$-(*). That is, we define a graph $H_\\Delta$ that satisfies $\\Delta$-(*), its chromatic number is $2\\Delta-2$ and every other graph that satisfies $\\Delta$-(*) admits a homomorphism in $H_\\Delta$.\nThe graph $H_{3}$ is depicted in \\cref{fig:h3}.\n\nRecall that the (categorical) product $G \\times H$ of graphs $G,H$ is the graph on $V(G) \\times V(H)$, such that $((g,h),(g',h')) \\in E(G \\times H)$ iff $(g,g') \\in E(G)$ and $(h,h') \\in E(H)$.\n \n Write $P$ for the product $K_{\\Delta-1}\\times K_{\\Delta-1}$. Let $V_0$ and $V_1$ be vertex disjoint copies of $K_{\\Delta-1}$.\nWe think of vertices in $V_i$ and $P$ as having labels from $[\\Delta-1]$ and $[\\Delta-1] \\times [\\Delta-1]$, respectively.\nThe graph $H_\\Delta$ is the disjoint union of $V_0$, $V_1$, $P$ and an extra vertex $\\dagger$ that is connected by an edge to every vertex in $P$, and additionally, if $v$ is a vertex in $V_0$ with label $i\\in [\\Delta-1]$, then we connect it by an edge with $(i',j)\\in P$ for every $i'\\not=i$ and $j\\in [\\Delta-1]$, and if $v$ is a vertex in $V_1$ with label $j\\in \\Delta-1$, then we connect it by an edge with $(i,j')\\in P$ for every $j'\\not=j$ and $i \\in [\\Delta-1]$.\n\n\n\\begin{proposition}\\label{pr:MinimalExample}\n\\begin{enumerate}\n \\item \\label{prop:hdelta}$H_\\Delta$ satisfies $\\Delta$-(*).\n \\item \\label{prop:hdeltachrom} $\\chi(H_\\Delta)=2\\Delta-2$.\n \\item \\label{prop:hdeltamax} A graph $G$ satisfies $\\Delta$-(*) if and only if it admits a homomorphism to $H_\\Delta$.\n\\end{enumerate}\n\n\\end{proposition}\n\\begin{proof}\n\\eqref{prop:hdelta} Set $S_0=V(V_0)\\cup \\{\\dagger\\}$ and $S_1=V(V_1)\\cup \\{\\dagger\\}$.\nBy the definition there are no edges between $S_0$ and $S_1$.\nConsider now, e.g., $V(H_\\Delta)\\setminus S_0$.\nLet $A_j$ consist of all elements in $P$ that have second coordinate equal to $j$ together with the vertex in $V_1$ that has the label $j$.\nBy the definition, the set $A_i$ is independent and $\\bigcup_{i\\in [\\Delta-1]}A_i$ covers $H_\\Delta\\setminus S_0$, and similarly for $S_1$.\n\n\\eqref{prop:hdeltachrom} By \\eqref{prop:hdelta} and the claim after the definition of $\\Delta$-(*), it is enough to show that $\\chi(H_\\Delta)\\geq 2\\Delta-2$. Towards a contradiction, assume that $c$ is a proper coloring of $H_\\Delta$ with $<2\\Delta-2$-many colors. Note the vertex $\\dagger$ guarantees that $|c(V(P))|\\leq 2\\Delta-4$, and also $\\Delta-1 \\leq |c(V(P))|$.\n\nFirst we claim that there are no indices $i,j\\in [\\Delta-1]$ (even with $i=j$) such that $c(i,r)\\not=c(i,s)$ and $c(r,j)\\not=c(s,j)$ for every $s\\not=r$: indeed, otherwise, by the definition of $P$ we had $c(i,r)\\not=c(s,j)$ for every $r,s$ unless $(i,r)=(s,j)$, which would the upper bound on the size of $c(V(P))$.\n\nTherefore, without loss of generality, we may assume that for every $i\\in [\\Delta-1]$ there is a color $\\alpha_i$ and two indices $j_i \\neq j'_i$ such that $c(i,j_i)=c(i,j'_i)=\\alpha_i$. It follows form the definition of $P$ and $j_i \\neq j'_i$ that\n $\\alpha_i\\not=\\alpha_{i'}$ whenever $i\\not= i'$. \n \n Moreover, note that any vertex in $V_1$ is connected to at least one of the vertices $(i,j_i)$ and $(i,j'_i)$, hence none of the colors $\\{\\alpha_i\\}_{i \\in [\\Delta-1]}$ can appear on $V_1$. Consequently, since $V_1$ is isomorphic to $K_{\\Delta-1}$ we need to use at least $\\Delta-1$ additional colors, a contradiction.\n\n\\eqref{prop:hdeltamax} First note that if $G$ admits a homomorphism into $H_\\Delta$, then the pullbacks of the sets witnessing $\\Delta$-(*) will witness that $G$ has $\\Delta$-(*).\n\nConversely, let $G$ be a graph that satisfies $\\Delta$-(*).\nFix the corresponding sets $S_0,S_1$ together with $(\\Delta-1)$-colorings $c_0,c_1$ of their complements.\nWe construct a homomorphism $\\Theta$ from $G$ to $H_\\Delta$. Let\n\\[\\Theta(v)=\n\\begin{cases} \\dagger &\\mbox{if } v \\in S_0 \\cap S_1, \\\\\nc_0(v) & \\mbox{if } v \\in S_1 \\setminus S_0, \\\\\nc_1(v) & \\mbox{if } v \\in S_0 \\setminus S_1,\n\\\\\n(c_0(v),c_1(v)) & \\mbox{if } v \\not \\in S_0 \\cup S_1.\n\\end{cases} \\]\n\n\nObserve that $S=S_0\\cap S_1$ is an independent set such that there is no edge between $S$ and $S_0 \\cup S_1$.\nUsing this observation, one easily checks case-by-case that $\\Theta$ is indeed a homomorphism.\n\n \n \n \n \n \n \n\\end{proof}\n\nNow, combining what we have so far, we can easily prove Theorem \\ref{th:homomorphism}.\n\n\n\\begin{proof}[Proof of Theorem \\ref{th:homomorphism}]\nIt follows that $\\Pi_{H_\\Delta}$ is not playable from \\cref{pr:MinimalExample}, \\cref{pr:CONdition*}. It is easy to see that if for a graph $G$ the LCL $\\Pi_G$ is not playable then $\\Pi_{G'}$ is not playable for every subgraph $G'$ of $G$. Since erasing a vertex decreases the chromatic number with at most one, for each $k\\leq 2\\Delta-2$ there is a subgraph $G_k$ of $H_\\Delta$ with $\\chi(G_k)=k$, such that $\\Pi_{G_k}$ is not playable.\n\nIt follows from Theorems \\ref{th:MainLOCAL} and \\ref{thm:MainBorel} that there is a $\\Delta$-regular Borel forest that admits no Borel homomorphism to any graph of $G_k$ and that $\\Pi_{G_k} \\not \\in \\mathsf{LOCAL}(O(\\log^*(n)))$.\n\nFinally, note that if $k \\geq \\Delta$ then $G_k$ can be chosen so that it contains $K_\\Delta$, yielding $\\Pi_{G_k} \\in \\mathsf{RLOCAL}(O(\\log\\log(n))$.\n\\end{proof}\n\n\\begin{remark}\n\\label{r:hedet}Recall that Hedetniemi's conjecture is the statement that if $G,H$ are finite graphs then $\\chi(G \\times H)=\\min\\{\\chi(G),\\chi(H)\\}$. This conjecture has been recently disproven by Shitov \\cite{shitov}, and strong counterexamples have been constructed later (see, \\cite{tardif2019note,zhu2021note}). We claim that these imply for $\\varepsilon>0$ the existence of finite graphs $H$ with $\\chi(H) \\geq (2-\\varepsilon)\\Delta$ to which $\\Delta$-regular Borel forests cannot have a homomorphism in BOREL, for every large enough $\\Delta$. Indeed, if a $\\Delta$-regular Borel forest admitted a Borel homomorphism to each finite graph of chromatic number at least $(2-\\varepsilon)\\Delta$, it would have such a homomorphism to their product as well. Thus, we would obtain that the chromatic number of the product of any graphs of chromatic number $(2-\\varepsilon)\\Delta$ is at least $\\Delta+1$. This contradicts Zhu's result \\cite{zhu2021note}, which states that the chromatic number of the product of graphs with chromatic number $n$ can drop to $\\approx \\frac{n}{2}$.\n\n\\end{remark}\n\n\\begin{remark}\nA natural attempt to construct graphs with large girth and not playable homomorphisms problem would be to consider random $d$-regular graphs of size $n$ for a large enough $n$. However, it is not hard to see that setting $\\Lambda(A)=\\operatorname{Alice}$ if and only if $|A|<\\frac{n}{d}$ shows that this approach cannot work.\n\\end{remark}\n\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{Proof of \\cref{th:MainLOCAL}}\\label{sec:marksproof}\n\nIn this section we prove Theorem \\ref{th:MainLOCAL}, by applying Marks' game technique in the LOCAL setting. In order to define our games, we will need certain auxiliary graphs, the so-called \\emph{ID graphs}. \nThe purpose of these graphs is to define a ``playground'' for the games that we consider.\nNamely, vertices in the game are labeled by vertices from the ID graph in such a way that the at the end we obtain a homomorphism from our underlying graph to the ID graph. \n\n\n\n\n\\begin{definition}\nA pair $\\idgraphcolored{n}{t}{r}=(\\idgraph{n}{t}{r},c)$ is called an \\emph{ID graph}, if\n\\begin{enumerate}\n \\item \\label{p:idgirth} $\\idgraph{n}{t}{r}$ is graph with girth at least $2t+2$,\n \\item \\label{p:idsize} $|V(\\idgraph{n}{t}{r})|\\le n$,\n \\item \\label{p:idlabels} $c$ is a $\\Delta$-edge-coloring of $\\idgraph{n}{t}{r}$, such that every vertex is adjacent to at least one edge of each color,\n \\item \\label{p:idratio} for each $\\alpha \\in [\\Delta]$ the ratio of a maximal independent set of $\\idgraph{n}{t}{r}^{\\alpha}=(V(\\idgraph{n}{t}{r}), E(\\idgraph{n}{t}{r}) \\cap c^{-1}(\\alpha))$ is at most $r$ (i.e., $\\idgraph{n}{t}{r}^\\alpha$ is the graph formed by $\\alpha$-colored edges).\n\\end{enumerate}\n\\end{definition}\n\nBefore we define the game we show that ID graphs exist.\n\n\\begin{proposition}\\label{pr:IDgraph}\nLet $t_n\\in o(\\log(n))$, $r>0$, $\\Delta\\geq 2$.\nThen there is an ID graph $\\idgraphcolored{n}{t_n}{r}$ for every $n\\in \\mathbb{N}$ sufficiently large.\n\\end{proposition}\n\\begin{proof}\nWe use the \\emph{configuration model} for regular random graphs, see \\cite{wormald1999models}.\nThis model is defined as follows.\nLet $n$ be even.\nThen a $d$-regular random sample on $n$-vertices is just a union of $d$-many independent uniform random perfect matchings.\nNote that in this model we allow parallel edges.\n\nIt was proved by Bollob\\'as \\cite{bollobas} that the independence ratio of a random $d$-regular graph is at most $\\frac{2\\log(d)}{d}$ a.a.s.\nMoreover, this quantity is concentrated by an easy application of McDiarmid's result \\cite{mcdiarmid2002concentration}, i.e., \n\\begin{equation}\n \\P\\left(\\left|X-\\textrm{\\textbf{E}}(X)\\right|\\ge \\sqrt{n}\\right)<2\\exp\\left(-\\frac{n}{d}\\right),\n\\end{equation}\nwhere $X$ is the random variable that counts the size of a maximal independent set.\nTherefore for fixed $n$ large enough we have that the independence ratio of a random sample is at most $\\frac{3\\log(d)}{d}$ with probability at least $\\left(1-2\\exp\\left(-\\frac{n}{d}\\right)\\right)$.\n\n\nPick a $d$ large enough such that $\\frac{3\\log(d)}{d} & 1-\\left(1-2\\exp\\left(-\\frac{n}{d}\\right)\\right)^\\Delta\n\\end{split}\n\\end{equation}\nas $n\\to \\infty$.\nThis shows that there exists such a graph $\\idgraphcolored{t_n}{n}{r}$ with non-zero probability.\n\\end{proof}\n\n\n\n\nNext, we define the games. As mentioned before, the games are going to depend on the following parameters: an algorithm $\\mathcal{A}_n$ of local complexity $t\\in o(\\log(n))$, an ID graph $\\idgraphcolored{n}{t}{r}$, $\\alpha\\in \\Delta$, $\\sigma\\in V(\\idgraph{n}{t}{r})$ and $S \\subseteq \\Sigma$.\n(We will view $\\mathcal{A}_n$, $\\idgraphcolored{n}{t}{r}$ as fixed, and the rest of the parameters as variables).\n\nThe game\n$$\\mathbb{G}(\\mathcal{A}_n,n,t,\\idgraphcolored{n}{t}{r})[\\alpha,\\sigma,S]$$\nis defined as follows: two players, $\\operatorname{Alice}$ and $\\operatorname{Bob}$ assign labels to the vertices of a rooted $\\Delta$-regular tree of diameter $t$. The labels are vertices of $\\idgraph{n}{t}{r}$ and the root is labeled $\\sigma$. In the $k$-th round, where $00$. $\\idgraphingcolored{r}=(\\idgraphing{r},c)$ is an \\emph{ID graphing}, if\n\n\\begin{enumerate}\n \\item $\\idgraphing{r}$ is an acyclic locally finite Borel graphing on a standard probability measure space $(X,\\mu)$,\n \\item \\label{p:idgraphinglabels} $c$ is a Borel $\\Delta$-edge-coloring of $\\idgraphing{r}$, such that every vertex is adjacent to at least one edge of each color,\n \\item \\label{p:idgraphingratio} for each $\\alpha \\in [\\Delta]$ the $\\mu$-measure of a maximal independent set of $\\idgraphing{r}^{\\alpha}=(V(\\idgraphing{r}), E(\\idgraphing{r}) \\cap c^{-1}(\\alpha))$ is at most $r$.\n\\end{enumerate}\n\\end{definition}\n\n\\begin{proposition}\n\\label{p:idgraphing}\nFor each $r>0$ there exists an ID graphing $\\idgraphingcolored{r}$.\n\\end{proposition}\n\\begin{proof}\nLet $\\idgraphing{r}$ be a \\emph{local-global limit} of the random graphs constructed in \\cref{pr:IDgraph} (see, e.g., \\cite{hatamilovaszszegedy} for the basic results about local-global convergence).\nIt is not hard to check that this limit satisfies the required properties.\n\\end{proof}\n\n\nNow we are ready to prove the theorem. The proof will closely follow the argument given in the proof of the LOCAL version, i.e., Theorem \\ref{th:MainLOCAL}, but can be understood without reading the latter.\n\n\\begin{proof}[Proof of \\cref{thm:MainBorel}]\nLet $\\Pi$ be an LCL that is not playable, and $N\\in \\mathbb{N}$ be the number of all possible colorings of vertices of $(\\mathbb{P}_{\\alpha})_{\\alpha\\in \\Delta}$ with two colors, that is, $N=2^{\\sum_\\alpha |V(\\mathbb{P}_\\alpha)|}$.\nSet $r:=\\frac{1}{N+1}$.\n\nWe define a Borel acyclic $\\Delta$-regular graph $\\mathcal{G}$, with edges properly colored by $\\Delta$, that does not admit a Borel solution of $\\Pi$.\nVertices of $\\mathcal{G}$ are pairs $(x,A)$, where $x\\in X$ is a vertex of $\\idgraphing{r}$ and $A$ is a countable subgraph of $\\idgraphing{r}$ that is a $\\Delta$-regular tree that contains $x$ and the edge coloring of $\\idgraphingcolored{r}$ induces a proper edge coloring of $A$.\nWe say that $(x,A)$ and $(y,B)$ are connected by an $\\alpha$-edge in $\\mathcal{G}$ if $A=B$, $x,y$ are adjacent in $A$ and the edge that connects them has color $\\alpha\\in \\Delta$.\n\nSuppose for a contradiction that $\\mathcal{A}$ is a Borel function that solves $\\Pi$ on $\\mathcal{G}$.\n\nNext, we define a family of games parametrized by $\\alpha \\in [\\Delta]$, $x \\in V(\\idgraphing{r})$ and $S \\subseteq \\Sigma$. For the reader familiar with Marks' construction, let us point out that for a fixed $x$, the games are analogues to the ones he defines, with the following differences: allowed moves are vertices of the ID graphing $\\idgraphing{r}$ and restricted by its edge relation, and the winning condition is defined by a set of labels, not just merely one label.\n\nSo, the game\n$$\\mathbb{G}(\\mathcal{A},\\idgraphingcolored{r})[\\alpha,x,S]$$ \nis defined as follows: $\\operatorname{Alice}$ and $\\operatorname{Bob}$ alternatingly label vertices of a $\\Delta$-regular rooted tree. The root is labelled by $x$, and the labels come from $V(\\idgraphing{r})$. In the $k$-th round, first $\\operatorname{Alice}$ labels vertices of distance $k$ from the root on the side of the $\\alpha$ edge.\nAfter that, $\\operatorname{Bob}$ labels all remaining vertices of distance $k$, etc (see \\cref{fig:game}). We also require the assignment of labels to give rise to an edge-color preserving homomorphism to $\\idgraphingcolored{r}$. \n\nIt follows from the acyclicity of $\\idgraphing{r}$ that a play of a game determines a $\\Delta$-regular rooted subtree of $\\idgraphing{r}$ to which the restriction of the edge-coloring is proper. That is, it determines a vertex $(x,A)$ of $\\mathcal{G}$. Let $\\operatorname{Alice}$ win iff the output of $\\mathcal{A}$ on the half-edge determined by $(x,A)$ and the color $\\alpha$ is \\textbf{not} in $S$.\n\nDefine the function $\\Lambda^x_{\\alpha}:\\mathbb{P}_\\alpha\\to \\{\\operatorname{Alice},\\operatorname{Bob}\\}$ assigning the player to some $S \\in V(\\mathbb{P}_\\alpha)$ who has a winning strategy in $\\mathbb{G}(\\mathcal{A},\\idgraphingcolored{r})[\\alpha,x,S]$.\nNote that, since $\\idgraphing{r}$ is locally finite, each player has only finitely many choices at each position. Thus, it follows from Borel Determinacy Theorem that $\\Lambda^x_{\\alpha}$ is well defined.\n\nNow we show the analogue of Proposition \\ref{pr:(A)satisfied}.\n\n\\begin{proposition}\\label{pr:Borel(A)satisfied}\n$(\\Lambda^{x}_{\\alpha})_{\\alpha\\in \\Delta}$ satisfies (A) in \\cref{def:Playability}.\n\\end{proposition}\n\\begin{proof}\nAssume that $(S_\\alpha)_{\\alpha\\in \\Delta}$ is such that $\\Lambda^x_\\alpha(S_\\alpha)=\\operatorname{Alice}$.\nThis means that $\\operatorname{Alice}$ has winning strategy in all the games corresponding to $S_\\alpha$. Letting these strategies play against each other in the obvious way, produces a vertex $(x,A)$ of $\\mathcal{G}$.\nSince $\\mathcal{A}$ solves $\\Pi$ it has to output labeling of half edges $(\\mathtt{a}_{\\alpha})_{\\alpha\\in \\Delta}\\in \\mathcal{N}$, where $\\mathtt{a}_\\alpha$ is a label on the half edge that start at $(x,A)$ and has color $\\alpha$.\nNote that we must have $\\mathtt{a}_\\alpha\\not\\in S_\\alpha$ by the definition of winning strategy for $\\operatorname{Alice}$.\nThis shows that (A) of Definition \\ref{def:Playability} holds.\n\\end{proof}\n\n Let $f$ be \n defined by $$x\\mapsto (\\Lambda^x_{\\alpha})_{\\alpha\\in\\Delta}.$$ Note that $f$ has a finite range. Using the fact that the allowed moves for each player can be determined in a Borel way, uniformly in $x$, it is not hard to see that for each element $s$ in the range, $f^{-1}(s)$ is in the algebra generated by sets that can be obtained by applying game quantifiers to Borel sets (for the definition see \\cite[Section 20.D]{kechrisclassical}). It has been shown by Solovay \\cite{Solovay} and independently by Fenstad-Norman \\cite{fenstadnorman} that such sets are provably $\\Delta^1_2$, and consequently, measurable. Therefore, $f$ is a measurable map.\n \n As the number of sequences of functions $(\\Lambda^x_{\\alpha})_{\\alpha\\in\\Delta}$ is $\\leq N$, by the choice of $r$, there exists a Borel set $Y$ with $\\mu(Y) > r$, such that $f$ is constant on $Y$.\n \nSince $\\Pi$ is not playable Proposition \\ref{pr:Borel(A)satisfied} gives that there is an $\\alpha\\in \\Delta$ and $S,T\\in \\mathbb{P}_\\alpha$ that violate condition (B).\nRecall that the measure of an independent Borel set of $\\idgraphing{r}^\\alpha$ is at most $r$.\nThat means that there are $x,y\\in Y$ such that $(x,y)$ is an $\\alpha$-edge in $\\idgraphing{r}$.\nWe let the winning strategies of $\\operatorname{Bob}$ in the games\n$$\\mathbb{G}(\\mathcal{A},\\idgraphingcolored{r})[\\alpha,x,S], \\ \\mathbb{G}(\\mathcal{A},\\idgraphingcolored{r})[\\alpha,y,T],$$\nplay against each other, where we start with an $\\alpha$ edge with endpoints labeled by $x,y$.\nThis produces a labeling of a $\\Delta$-regular tree $A$ with labels from $\\idgraphing{r}$ such that $(x,A)$ and $(y,A)$ span an $\\alpha$-edge in $\\idgraphing{r}$.\nApplying $\\mathcal{A}$ on the half-edge determined by $(x,A)$ and the color $\\alpha$, we obtain $\\mathtt{a}_0\\in S$.\nSimilarly, we produce $\\mathtt{a}_1\\in T$ for $(y,A)$.\nHowever, $(\\mathtt{a}_0,\\mathtt{a}_1)\\not\\in {\\bf E}_\\alpha$ by the definition of an edge in $\\mathbb{P}_{\\alpha}$.\nThis shows that $\\mathcal{A}$ does not produce a Borel solution to $\\Pi$.\n\\end{proof}\n\n\\begin{remark}\nOne can give an alternative proof of the existence of a Borel $\\Delta$-regular forest that does not admit a Borel homomorphism to $G_k$ that avoids using the graphing $\\idgraphing{r}$ and uses the original example described by Marks \\cite{DetMarks} instead.\nThe reason is that the example graphs $G_k$ contain $K_\\Delta$, as discussed in the proof of \\cref{th:homomorphism}.\nThis takes care of the ``non-free'' as in the argument of Marks.\nThen it is enough to use the pigeonhole principle on $\\mathbb{N}$ instead of the independence ratio reasoning.\n\\end{remark}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Separation Of Various Complexity Classes}\n\\label{sec:separating_examples}\n\nIn this section we provide examples of several problems that separate some classes from \\cref{fig:big_picture_trees}.\nThe examples show two things. First, we have $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon) \\not= \\mathsf{OLOCAL}(O(\\log\\log 1\/\\varepsilon))$. \nThis shows that there are problems such that their worst case complexity is at least $\\Theta(\\log n)$ on finite $\\Delta$-regular trees, but their average local complexity is constant. \nSecond, we show that there are problems in the class $\\mathsf{BOREL}$ that are not in the class $\\mathsf{RLOCAL}(O(\\log\\log n)) = \\mathsf{RLOCAL}(o(\\log n))$, on $\\Delta$-regular trees. This shows that one cannot in general hope that results from (Borel) measurable combinatorics can be turned into very efficient (sublogarithmic) distributed algorithms. \n\n\\subsection{Preliminaries}\n\\label{subsec:preliminaries_separating}\n\nWe first show that there is a strong connection between randomized local complexities and uniform local complexities. Afterwards, we introduce a generic construction that turns any LCL into a new LCL. We later use this generic transformation to construct an LCL that is contained in the set $\\mathsf{BOREL} \\setminus \\mathsf{RLOCAL}(O(\\log \\log n))$.\n\n\\paragraph{Uniform vs Local Randomized Complexity} \n\nWe will now discuss the connections between uniform and randomized complexities. \nNote that the easy part of the connection between the two concepts is turning uniform local algorithms into randomized ones, as formalized in the following proposition. \n\n\\begin{proposition}\n\\label{prop:bee_gees}\nWe have $\\mathsf{OLOCAL}(t(\\varepsilon)) \\subseteq \\mathsf{RLOCAL}(t(1\/n^{O(1)}))$. \n\\end{proposition}\n\\begin{proof}\nWe claim that an uniform local algorithm $\\mathcal{A}$ with an uniform local complexity of $t(\\varepsilon)$ can be turned into a local randomized algorithm $\\mathcal{A}'$ with a local complexity of $t(1\/n^{O(1)})$. \n\nThe algorithm $\\mathcal{A}'$ simulates $\\mathcal{A}$ on an infinite $\\Delta$-regular tree -- each vertex $u$ of degree less than $\\Delta$ in the original tree pretends that the tree continues past its virtual half-edges and the random bits of $u$ are used to simulate the random bits in this virtual subtree. \nChoosing $\\varepsilon = 1\/n^{O(1)}$, one gets that the probability of $\\mathcal{A}'$ needing to look further than $t(\\varepsilon)$ for any vertex is bounded by $1\/n^{O(1)}$, as needed in the definition of the randomized local complexity. \n\\end{proof}\n\n\nOn the other hand, we will use the following proposition from \\cite{grebik_rozhon2021toasts_and_tails}. It informally states that the existence of \\emph{any} uniform local algorithm together with the existence of a sufficiently fast randomized local algorithm for a given LCL $\\Pi$ directly implies an upper bound on the uniform complexity of $\\Pi$.\n\\begin{proposition}[\\cite{grebik_rozhon2021toasts_and_tails}]\n\\label{prop:local->uniform}\nLet $\\mathcal{A}$ be an uniform local algorithm solving an LCL problem $\\Pi$ such that its \\emph{randomized} local complexity on finite $\\Delta$-regular trees is $t(n)$ for $t(n) = o(\\log n)$. \nThen, the uniform local complexity of $\\mathcal{A}$ on infinite $\\Delta$-regular trees is $O(t(1\/\\varepsilon))$.\n\\end{proposition}\n\n\\cref{prop:local->uniform} makes our life simpler, since we can take a known randomized distributed local algorithm with local complexity $g(n) = o(\\log n)$ and check if it works without the knowledge of $n$. If yes, this automatically implies that the uniform local complexity of the algorithm is $h(\\varepsilon) = O(g(1\/\\varepsilon))$. \nIn particular, combining \\cref{prop:local->uniform} with the work of \\cite{chang_kopelowitz_pettie2019exp_separation,chang_pettie2019time_hierarchy_trees_rand_speedup} and verifying that the algorithms of \nFischer and Ghaffari~\\cite[Section 3.1.1]{fischer-ghaffari-lll} and Chang~et~al.~\\cite[Section 5.1]{ChangHLPU20} work without the knowledge of $n$, we obtain the following result. \n\n\n\\begin{theorem}\n\\label{thm:classification_local_uniform_trees}\nWe have:\n\\begin{itemize}\n \\item $\\mathsf{LOCAL}(O(1)) = \\mathsf{OLOCAL}(O(1))$\n \\item $\\mathsf{RLOCAL}(O(\\log^* n)) = \\mathsf{OLOCAL}(O(\\log^* 1\/\\varepsilon))$\n \\item $\\mathsf{RLOCAL}(O(\\log\\log n)) = \\mathsf{OLOCAL}(O(\\log\\log 1\/\\varepsilon))$\n\\end{itemize}\nMoreover, there are no other possible uniform local complexities for $t(\\varepsilon) = o(\\log 1\/\\varepsilon)$. \n\\end{theorem}\nWe note that the first two items are proven in \\cite{grebik_rozhon2021toasts_and_tails}.\nFor the third item it suffices by known reductions to find an uniform algorithm\nsolving a version of the distributed Lov\\'asz Local Lemma (LLL) on so-called tree-structured dependency graphs considered in~\\cite{ChangHLPU20}.\n\n\\begin{proof}[Proof sketch]\nThe proof follows from \\cref{prop:bee_gees} and the following ideas. \nThe first item follows from the fact that any local algorithm with local complexity $O(1)$ can simply be made uniform. More specifically, there exists a constant $n_0$ --- depending only on the problem $\\Pi$ --- such that the algorithm, being told the size of the graph is $n_0$, is correct on \\emph{any} graph of size $n \\ge n_0$. \n\nFor the second item, by the work of \\cite{chang_kopelowitz_pettie2019exp_separation,chang_pettie2019time_hierarchy_trees_rand_speedup}, it suffices to check that there is an uniform distributed $(\\Delta + 1)$-coloring algorithm with uniform local complexity $O(\\log^* 1\/\\varepsilon)$. Such an algorithm was given in \\cite{HolroydSchrammWilson2017FinitaryColoring}, or follows from the work in \\cite{Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring} and \\cref{prop:local->uniform}. \n\nSimilarly, for the third item, by the work of \\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} it suffices to check that there is an uniform distributed algorithm for a specific LLL problem on trees\nwith uniform local complexity $O(\\log\\log 1\/\\varepsilon)$.\nSuch an algorithm can be obtained by combining the randomized \\emph{pre-shattering} algorithm of Fischer and Ghaffari~\\cite[Section 3.1.1]{fischer-ghaffari-lll} and the deterministic \\emph{post-shattering} algorithm of Chang~et~al.~\\cite[Section 5.1]{ChangHLPU20} in a \\emph{graph shattering} framework~\\cite{barenboim2016locality}, which solves the LLL problem \nwith local complexity $O(\\log\\log n)$.\nBy \\cref{prop:local->uniform}, it suffices to check that this algorithm can be made to work even if it does not know the size of the graph $n$. We defer the details to \\cref{sec:LLL}.\nThis finishes the proof of \\cref{thm:classification_local_uniform_trees}. \n\\end{proof}\n\n\n\n\n\n\n\n\\paragraph{Adding Paths}\nBefore we proceed to show some separation results, we define a certain construction that turns any LCL problem $\\Pi$ into a new LCL problem $\\overline{\\Pi}$ with the following property. If the original problem $\\Pi$ cannot be solved by a fast local algorithm, then the same holds for $\\overline{\\Pi}$. However, $\\overline{\\Pi}$ might be strictly easier to solve than $\\Pi$ for $\\mathsf{BOREL}$ constructions. \n\n\\begin{definition}\nLet $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ be an LCL. \nWe define an LCL $\\overline{\\Pi}=(\\Sigma',\\mathcal{V}',\\mathcal{E}')$ as follows.\nLet $\\Sigma'$ be $\\Sigma$ together with one new label.\nLet $\\mathcal{V}'$ be the union of $\\mathcal{V}$ together with any cardinality-$\\Delta$ multiset that contains the new label exactly two times.\nLet $\\mathcal{E}'$ be the union of $\\mathcal{E}$ together with the cardinality-$2$ multiset that contains the new label twice.\n\\end{definition}\n\nIn other words, the new label determines doubly infinite lines in infinite $\\Delta$-regular trees, or lines that start and end in virtual half-edges in finite $\\Delta$-regular trees.\nMoreover, a vertex that is on such a line does not have to satisfy any other vertex constraint.\nWe call these vertices \\emph{line-vertices} and each edge on a line a \\emph{line-edge}.\n\n\n\\begin{proposition}\n\\label{pr:AddingLineLLL}\nLet $\\Pi$ be an LCL problem such that $\\overline{\\Pi}\\in \\mathsf{RLOCAL}(t(n))$ for $t(n) = o(\\log(n))$. Then also $\\Pi \\in \\mathsf{RLOCAL}(O(t(n)))$.\n\\end{proposition}\n\\begin{proof}\nLet $\\overline{\\mathcal{A}}$ be a randomized $\\mathsf{LOCAL}$ algorithm that solves $\\overline{\\Pi}$ in $t(n) = o(\\log n)$ rounds with probability at least $1 - 1\/n^C$ for some sufficiently large constant $C$. We will construct an algorithm $\\mathcal{A}$ for $\\Pi$ with complexity $t(n)$ that is correct with probability $1 - \\frac{4}{n^{C\/3 - 2}}$. The success probability of $\\mathcal{A}$ can then be boosted by ``lying to it'' that the number of vertices is $n^{O(1)}$ instead of $n$; this increases the running time by at most a constant factor and boosts the success probability back to $1 - 1\/n^C$. \n\nConsider a $\\Delta$-regular rooted finite tree $T$ of depth $10t(n)$ and let $u$ be its root vertex. \nNote that $|T| \\le \\Delta^{10t(n)+1} < n$, for $n$ large enough. \n\n\nWe start by proving that when running $\\overline{\\mathcal{A}}$ on the tree $T$, then $u$ is most likely not a line-vertex. This observation then allows us to turn $\\overline{\\mathcal{A}}$ into an algorithm $\\mathcal{A}$ that solves $\\Pi$ on any $\\Delta$-regular input tree.\nLet $X$ be the indicator of $\\mathcal{A}$ marking $u$ as a line-vertex. Moreover, for $i \\in [\\Delta]$, let $Y_i$ be the indicator variable for the following event. The number of line edges in the $i$-th subtree of $u$ with one endpoint at depth $5t(n)$ and the other endpoint at depth $5t(n) + 1$ is odd.\n\nBy a simple parity argument we can relate the value of $X$ with the values of $Y_1,Y_2,\\ldots,Y_\\Delta$ as follows. If $X = 0$, that is, $u$ is not a line-vertex, then all of the $Y_i$'s have to be $0$, as each path in the tree $T$ is completely contained in one of the $\\Delta$ subtrees of $u$. On the other hand, if $u$ is a line-vertex, then there exists exactly one path, the one containing $u$, that is not completely contained in one of the $\\Delta$ subtrees of $u$. This in turn implies that exactly two of the $\\Delta$ variables $Y_1,Y_2,\\ldots,Y_\\Delta$ are equal to $1$.\n\nThe random variables $Y_1,Y_2,\\ldots,Y_\\Delta$ are identically distributed and mutually independent. Hence, if $\\P(Y_i = 1) > \\frac{1}{n^{C\/3}}$, then the probability that there are at least $3$ $Y_i$'s equal to $1$ is strictly greater than $\\left(\\frac{1}{n^{C\/3}}\\right)^3 = \\frac{1}{n^C}$. This is a contradiction, as in that case $\\overline{\\mathcal{A}}$ does not produce a valid output, which according to our assumption happens with probability at most $\\frac{1}{n^C}$. \n\nThus, we can conclude that $\\P(Y_i = 1) \\leq \\frac{1}{n^{C\/3}}$. By a union bound, this implies that all of the $Y_i$'s are zero with probability at least $\\frac{1}{n^{C\/3 - 1}}$, and in that case $u$ is not a line-vertex. \n\nFinally, the algorithm $\\mathcal{A}$ simulates $\\overline{\\mathcal{A}}$ as if the neighborhood of each vertex was a $\\Delta$-regular branching tree up to depth at least $t(n)$. \n\nIt remains to analyze the probability that $\\mathcal{A}$ produces a valid solution for the LCL problem $\\Pi$. To that end, let $v$ denote an arbitrary vertex of the input tree. The probability that the output of $v$'s half edges satisfy the vertex constraint of the LCL problem $\\overline{\\Pi}$ is at least $1 - 1\/n^C$. Moreover, in the case that $v$ is not a line-vertex, which happens with probability at least $1 - \\frac{1}{n^{C\/3 - 1}}$, the output of $v$'s half edges even satisfy the vertex constraint of the LCL problem $\\Pi$. Hence, by a union bound the vertex constraint of the LCL problem $\\Pi$ around $v$ is satisfied with probability at least $1 - \\frac{2}{n^{C\/3 - 1}}$. With exactly the same reasoning, one can argue that the probability that the output of $\\mathcal{A}$ satisfies the edge constraint of the LCL problem $\\Pi$ at a given edge is also at least $1 - \\frac{2}{n^{C\/3 - 1}}$. Finally, doing a union bound over the $n$ vertex constraints and $n-1$ edge constraints it follows that $\\mathcal{A}$ produces a valid solution for $\\Pi$ with probability at least $1 - \\frac{4}{n^{C\/3 - 2}}$.\n\\end{proof}\n\n\\begin{remark}\nThe ultimate way how to solve LCLs of the form $\\overline{\\Pi}$ is to find a spanning forest of lines.\nThis is because once we have a spanning forest of lines, then we might declare every vertex to be a line-vertex and label every half-edge that is contained on a line with the new symbol in $\\overline{\\Pi}$.\nThe remaining half-edges can be labeled arbitrarily in such a way that the edge constraints are satisfied.\n\nPut otherwise, once we are able to construct a spanning forest of lines and $\\mathcal{E}\\not=\\emptyset$, where $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$, then $\\overline{\\Pi}$ can be solved.\nMoreover, in that case the complexity of $\\overline{\\Pi}$ is upper bounded by the complexity of finding the spanning forest of lines.\n\nLyons and Nazarov \\cite{lyons2011perfect} showed that $\\Pi_{\\text{pm}}\\in \\mathsf{fiid}$ and it is discussed in \\cite{LyonsTrees} that the construction can be iterated $(\\Delta-2)$-many times. After removing each edge that is contained in one of the $\\Delta-2$ perfect matchings each vertex has a degree of two.\nHence, it is possible to construct a spanning forest of lines as fiid.\nConsequently, $\\overline{\\Pi}\\in \\mathsf{fiid}$ for every $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ with $\\mathcal{E}\\not=\\emptyset$.\n\\end{remark}\n\n\n\\subsection{$\\mathsf{LOCAL}(O(\\log^* n)) \\not= \\mathsf{BOREL}$}\n\nWe now formalize the proof that $\\mathsf{LOCAL}(O(\\log^* n)) \\not= \\mathsf{BOREL}$ from \\cref{sec:introseparation}. \nRecall that $\\Pi_{\\chi,\\Delta}$ is the proper vertex $\\Delta$-coloring problem.\nThe following claim directly follows by \\cref{pr:AddingLineLLL} together with the fact that $\\Pi_{\\chi,\\Delta}\\not\\in \\mathsf{LOCAL}(O(\\log^* n))$.\n\n\\begin{claim}\nWe have $\\overline{\\Pi_{\\chi,\\Delta}}\\not \\in \\mathsf{LOCAL}(O(\\log^* n))$.\n\\end{claim}\n\nOn the other hand we show that adding lines helps to find a Borel solution.\nThis already provides a simple example of a problem in the set $\\mathsf{BOREL} \\setminus \\mathsf{LOCAL}(O(\\log^* n))$.\n\n\\begin{proposition}\nWe have $\\overline{\\Pi_{\\chi,\\Delta}} \\in \\mathsf{BOREL}$. \n\\end{proposition}\n\\begin{proof}\nBy \\cite{KST}, every Borel graph of finite maximum degree admits a Borel maximal independent set.\nLet $\\mathcal{G}$ be a Borel $\\Delta$-regular acyclic graph on a standard Borel space $X$.\nBy an iterative application of the fact above we find Borel sets $A_1,\\dots,A_{\\Delta-2}$ such that $A_1$ is a maximal independent set in $\\mathcal{G}$ and $A_i$ is a maximal independent set in $\\mathcal{G}\\setminus \\bigcup_{j1$.\nWe think of $A_1,\\dots,A_{\\Delta-2}$ as color classes for the first $\\Delta-2$ colors.\nLet $B=X\\setminus \\bigcup_{i\\in [\\Delta-2]}A_i$ and let $\\mathcal{H}$ be the subgraph of $\\mathcal{G}$ determined by $B$.\nIt is easy to see that the maximum degree of $\\mathcal{H}$ is $2$.\nIn particular, $\\mathcal{H}$ consists of finite paths, one-ended infinite paths or doubly infinite paths.\nWe use the extra label in $\\overline{\\Pi_{\\chi,\\Delta}}$ to mark the doubly infinite paths.\nIt remains to use the remaining $2$ colors in $[\\Delta]$ to define a proper vertex $2$-coloring of the finite and one-ended paths.\nIt is a standard argument that this can be done in a Borel way and it is easy to verify that, altogether, we found a Borel $\\overline{\\Pi_{\\chi,\\Delta}}$-coloring of $\\mathcal{G}$.\n\\end{proof}\n\n\n\n\\subsection{Examples and Lower Bound}\n\nIn this subsection we define two LCLs and show that they are not in the class $\\mathsf{RLOCAL}(o(\\log n))$.\nIn the next subsection we show that one of them is in the class $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon)$ and the other in the class $\\mathsf{BOREL}$.\nBoth examples that we present are based on the following relaxation of the perfect matching problem.\n\n\\begin{definition}[Perfect matching in power-$2$ graph]\nLet $\\Pi^2_{\\text{pm}}$ be the perfect matching problem in the power-$2$ graph, i.e., in the graph that we obtain from the input graph by adding an edge between any two vertices that have distance at most $2$ in the input graph\n\\end{definition}\n\nWe show that $\\Pi^2_{\\text{pm}}$ is not contained in $\\mathsf{RLOCAL}(o(\\log n))$. \\cref{pr:AddingLineLLL} then directly implies that $\\overline{\\Pi^2_{\\text{pm}}}$ is not contained in $\\mathsf{RLOCAL}(o(\\log n))$ as well.\nOn the other hand, we later use a one-ended spanning forest decomposition to show that $\\Pi^2_{\\text{pm}}$ is in $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon)$ and a one or two-ended spanning forest decomposition to show that $\\overline{\\Pi^2_{\\text{pm}}}$ is in $\\mathsf{BOREL}$.\n\n\nWe first show the lower bound result.\nThe proof is based on a simple parity argument as in the proof of \\cref{pr:AddingLineLLL}.\n\n\\begin{theorem}\nThe problems $\\Pi^2_{\\text{pm}}$ and $\\overline{\\Pi^2_{\\text{pm}}}$ are not in the class $\\mathsf{RLOCAL}(o(\\log n))$.\n\\end{theorem}\n\n\\begin{proof}\nSuppose that the theorem statement does not hold for $\\Pi^2_{\\text{pm}}$.\nThen, by \\cref{thm:basicLOCAL}, there is a distributed randomized algorithm solving $\\Pi^2_{\\text{pm}}$ with probability at least $1 - 1\/n$. By telling the algorithm that the size of the input graph is $n^{2\\Delta}$ instead of $n$, we can further boost the success probability to $1 - 1\/n^{2\\Delta}$. The resulting round complexity is at most a constant factor larger, as $\\Delta = O(1)$.\nWe can furthermore assume that the resulting algorithm $\\mathcal{A}$ will always output for each vertex $v$ a vertex $M(v) \\neq v$ in $v$'s $2$-hop neighborhood, even if $\\mathcal{A}$ fails to produce a valid solution. The vertex $M(v)$ is the vertex $v$ decides to be matched to, and if $\\mathcal{A}$ produces a valid solution it holds that $M(M(v)) = v$.\n\nConsider a fully branching $\\Delta$-regular tree $T$ of depth $10t(n)$. Note that $|T| < n$ for $n$ large enough. Let $u$ be the root of $T$ and $T_i$ denote the $i$-th subtree of $u$ (so that $u \\not\\in T_i$).\nLet $S_i$ denote the set of vertices $v$ in $T_i$ such that both $v$ and $M(v)$ have distance at most $2t(n)$ from $u$.\nNote that $M(v)$ is not necessarily a vertex of $T_i$.\nLet $Y_i$ be the indicator of whether $S_i$ is even.\n\nObserve that, if $\\mathcal{A}$ does not fail on the given input graph, then the definition of the $S_i$ implies that every vertex in $\\{ u \\} \\cup \\bigcup_{i \\in [\\Delta]} S_i$ is matched to a vertex in $\\{ u \\} \\cup \\bigcup_{i \\in [\\Delta]} S_i$.\nHence, the number of vertices in $\\{ u \\} \\cup \\bigcup_{i \\in [\\Delta]} S_i$ is even, which implies that we cannot have $Y_1 = Y_2 = \\ldots = Y_\\Delta = 1$ unless $\\mathcal{A}$ fails.\nNote that $Y_i$ depends only on the output of $\\mathcal{A}$ at the vertices in $T_i$ that have a distance of precisely $2t(n) - 1$ or $2t(n)$ to $u$ (as all other vertices in $T_i$ are guaranteed to be in $S_i$). \nHence, the events $Y_i = 1$, $i \\in [\\Delta]$, are independent, and, since $\\P(Y_1 = 1) = \\ldots = \\P(Y_{\\Delta} = 1)$ (as all vertices in all $T_i$ see the same topology in their $t(n)$-hop view), we obtain $(\\P(Y_1 = 1))^{\\Delta} \\le 1\/n^{2\\Delta}$, which implies $\\P(Y_i = 1) \\le 1\/n^2$, for any $i \\in [\\Delta]$.\n\n\nHence, with probability at least $1 - \\Delta \/n^2$ we have $Y_1=Y_2=\\ldots = Y_\\Delta = 0$, by a union bound.\nLet $v_1, v_2, \\dots, v_\\Delta$ denote the neighbors of $u$ with $v_i$ being the root of subtree $T_i$.\nNote that if $\\mathcal{A}$ does not fail and $Y_1=Y_2=\\ldots = Y_\\Delta = 0$, then it has to be the case that $u$ is matched to a vertex in some subtree $T_i$ (not necessarily $v_i$), while any vertex from $N_{-i}(u) := \\{v_1, v_2, \\dots, v_{i-1}, v_{i+1}, \\dots, v_\\Delta\\}$ is matched with another vertex from $N_{-i}(u)$ (hence, we already get that $\\Delta$ needs to be odd).\nThis is because each subtree needs to have at least one vertex matched to a different subtree (since $|S_i|$ is odd for each $i \\in [\\Delta]$).\nIf $M(u)$ is a vertex in $T_i$ and, for any $v \\in N_{-i}(u)$, we have $M(v) \\in N_{-i}(u)$, we say that $u$ orients the edge going to $v_i$ outwards.\n\nConsider a path $(u_0 = u, u_1, \\dots, u_k)$, where $k = 2t(n) + 3$.\nBy the argument provided above, the probability that $u$ orients an edge outwards is at least $1 - \\Delta\/n^2 - 1\/n^{2\\Delta}$, and, if $u$ orients an edge outwards, the probability that this edge is not $uu_1$ is $(\\Delta - 1)\/\\Delta$, by symmetry.\nHence, we obtain (for sufficiently large $n$) that $\\P(\\mathcal{E}_1) \\geq 1\/2$, where $\\mathcal{E}_1$ denotes the event that vertex $u$ orients an edge different from $uu_1$ outwards.\nWith an analogous argument, we obtain that $\\P(\\mathcal{E}_2) \\geq 1\/2$, where $\\mathcal{E}_2$ denotes the event that vertex $u_k$ orients an edge different from $u_ku_{k-1}$ outwards.\nSince $\\mathcal{E}_1$ and $\\mathcal{E}_2$ are independent (due to the fact that the distance between any neighbor of $u$ and any neighbor of $u_k$ is at least $2t(n) + 1$), we obtain $\\P(\\mathcal{E}_1 \\cap \\mathcal{E}_2) \\geq 1\/4$.\nMoreover, the probability that all vertices in $\\{ u_1, \\dots, u_{k-1} \\}$ orient an edge outwards is at least $1 - (2t(n)+2)(\\Delta\/n^2 + 1\/n^{2\\Delta})$, which is at least $4\/5$ for sufficiently large $n$.\nThus, by a union bound, there is a probability of at least $1 - 3\/4 - 1\/5 - 1\/n^{2\\Delta}$ that $\\mathcal{A}$ does not fail, all vertices in $(u, u_1, \\dots, u_k)$ orient an edge outwards, and the outward oriented edges chosen by $u$ and $u_k$ are not $uu_1$ and $u_ku_{k-1}$, respectively.\nFor sufficiently large $n$, the indicated probability is strictly larger than $0$.\nWe will obtain a contradiction and conclude the proof for $\\Pi^2_{\\text{pm}}$ by showing that there is no correct solution with the described properties.\n\nAssume in the following that the solution provided by $\\mathcal{A}$ is correct and $u, u_1, \\dots, u_k$ satisfy the mentioned properties.\nSince $u$ orients an edge different from $uu_1$ outwards, $u_1$ must be matched to a neighbor of $u$, which implies that $u_1$ orients edge $u_1u$ outwards.\nUsing an analogous argumentation, we obtain inductively that $u_j$ orients edge $u_ju_{j-1}$ outwards, for all $2 \\leq j \\leq k$.\nHowever, this yields a contradiction to the fact that the edge that $u_k$ orients outwards is different from $u_ku_{k-1}$, concluding the proof (for $\\Pi^2_{\\text{pm}}$).\n\nThe theorem statement for $\\overline{\\Pi^2_{\\text{pm}}}$ follows by applying \\cref{pr:AddingLineLLL}.\n\n\n\\end{proof}\n\n\n\n\n\n\\subsection{Upper Bounds Using Forest Decompositions}\n\nWe prove an upper bound for the problems $\\Pi^2_{\\text{pm}}$ and $\\overline{\\Pi^2_{\\text{pm}}}$ defined in the previous subsection.\nWe use a technique of decomposing the input graph in a spanning forest with some additional properties, i.e., one or two ended. This technique was used in \\cite{BrooksMeas} to prove Brooks' theorem in a measurable context.\nNamely, they proved that if $\\mathcal{G}$ is a Borel $\\Delta$-regular acyclic graph and $\\mu$ is a Borel probability measure, then it is possible to erase some edges of $\\mathcal{G}$ in such a way that the remaining graph is a one-ended spanning forest on a $\\mu$-conull set.\nIt is not hard to see that one can solve $\\Pi^2_{\\text{pm}}$ on one-ended trees in an inductive manner, starting with the leaf vertices. Consequently we have $\\Pi^2_{\\text{pm}}\\in \\mathsf{MEASURE}$.\nWe provide a quantitative version of this result and thereby show that the problem of constructing a one-ended forest is contained in $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon)$.\nA variation of the construction shows that it is possible to find a one or two-ended spanning forest in a Borel way.\nThis allows us to show that $\\overline{\\Pi^2_{\\text{pm}}}\\in \\mathsf{BOREL}$.\nAs an application of the decomposition technique we show that Vizing's theorem, i.e., proper $\\Delta+1$ edge coloring $\\Pi_{\\chi',\\Delta+1}$, for $\\Delta=3$ is in the class $\\mathsf{OLOCAL}(\\operatorname{poly} \\log 1\/\\varepsilon)$.\n\n\n\\subsubsection{Uniform Complexity of the One-Forest Decomposition }\nWe start with the definition of a one-ended spanning forest.\nIt can be viewed as a variation of the edge grabbing problem $\\Pi_{\\text{edgegrab}}$.\nNamely, if a vertex $v$ grabs an incident edge $e$, then we think of $e$ as being oriented away from $v$ and $v$ selecting the other endpoint of $e$ as its parent.\nSuppose that $\\mathcal{T}$ is a solution of $\\Pi_{\\text{edgegrab}}$ on $T_\\Delta$.\nWe denote with $\\mathcal{T}^{\\leftarrow}(v)$ the subtree with root $v$.\n\n\n\\begin{definition}[One-ended spanning forest]\nWe say that a solution $\\mathcal{T}$ of $\\Pi_{\\text{edgegrab}}$ is a \\emph{one-ended spanning forest} if $\\mathcal{T}^{\\leftarrow}(v)$ is finite for every $v\\in T_\\Delta$ (see \\cref{fig:one_ended}). \n\n\\end{definition}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=.4\\textwidth]{fig\/one_ended.pdf}\n \\caption{An example of a one-ended forest decomposition of an infinite 3-regular tree. }\n \\label{fig:one_ended}\n\\end{figure}\n\n\nNote that on an infinite $\\Delta$-regular tree every connected component of a one-ended spanning forest must be infinite.\nFurthermore, recall that we discussed above that it is possible to construct a one-ended spanning forest in $\\mathsf{MEASURE}$ by \\cite{BrooksMeas}.\nNext we formulate our quantitative version.\n\n\n\n\n\n\\begin{restatable}{theorem}{oneendedtreetail}\n\\label{thm:one_ended_forest}\nThe one-ended spanning forest can be constructed by an uniform local algorithm with an uniform local complexity of $O(\\operatorname{poly}\\log1\/\\varepsilon)$. \nMore precisely, there is an uniform distributed algorithm $\\mathcal{A}$ that computes a one-ended spanning forest and if we define $R(v)$ to be the smallest coding radius that allows $v$ to compute $\\mathcal{T}^{\\leftarrow}(v)$, then\n\\[\nP(R(v) > O(\\operatorname{poly} \\log 1\/\\varepsilon)) \\le \\varepsilon. \n\\]\n\\end{restatable}\n\nThe formal proof of \\cref{thm:one_ended_forest} can be found in \\cref{sec:formalproof}.\nWe first show the main applications of the result.\nNamely, we show that $\\Pi^2_{\\text{pm}}$ can be solved inductively starting from the leaves of the one-ended trees. This proves that $\\Pi^2_{\\text{pm}}$ is in $\\mathsf{OLOCAL}(\\operatorname{poly} \\log 1\/\\varepsilon)$ and hence $\\mathsf{OLOCAL}(\\operatorname{poly} \\log 1\/\\varepsilon) \\setminus \\mathsf{RLOCAL}(o(\\log n))$ is non-empty.\n\n\\begin{theorem}\\label{thm:twomatchFFIID}\nThe problem $\\Pi^2_{\\text{pm}}$ is in the class $\\mathsf{MEASURE}$ and $\\mathsf{OLOCAL}(O(\\operatorname{poly} \\log 1\/\\varepsilon))$.\n\\end{theorem}\n\\begin{proof}\nFirst we show that every infinite one-ended directed tree $T$ admits an inductive solution from the leaves to the direction of infinity. More precisely, we define a power-$2$ perfect matching on the infinite one-ended directed tree $T$ such that the following holds. Let $v$ be an arbitrary vertex and $T^{\\leftarrow}(v)$ the finite subtree rooted at $v$. Then, all vertices in $T^{\\leftarrow}(v)$, with the possible exception of $v$, are matched to a vertex in $T^{\\leftarrow}(v)$. Moreover, for each vertex in $T^{\\leftarrow}(v) \\setminus \\{v\\}$, it is possible to determine to which vertex it is matched by only considering the subtree $T^{\\leftarrow}(v)$. \nWe show that it is possible to define a perfect matching in that way by induction on the height of the subtree $T^{\\leftarrow}(v)$. \n\nWe first start with the leaves and declare them as unmatched in the first step of the induction.\n\nNow, consider an arbitrary vertex $v$ with $R \\geq 1$ children that we denote by $v_1, v_2, \\ldots, v_R$.\nBy the induction hypothesis, all vertices in $T^{\\leftarrow}(v_i) \\setminus \\{v_i\\}$ are matched with a vertex in $T^{\\leftarrow}(v_i)$ for $i \\in [R]$ and it is possible to compute that matching given $T^{\\leftarrow}(v)$. However, some of the children of $v$ are possibly still unmatched. If the number of unmatched children of $v$ is even, then we simply pair them up in an arbitrary but fixed way and we leave $v$ unmatched. Otherwise, if the number of unmatched children of $v$ is odd, we pair up $v$ with one of its unmatched children and we pair up all the remaining unmatched children of $v$ in an arbitrary but fixed way. This construction clearly satisfies the criteria stated above. In particular, the resulting matching is a perfect matching, as every vertex gets matched eventually.\n\n\nBy \\cref{thm:one_ended_forest} there is an uniform distributed algorithm of complexity $O(\\operatorname{poly} \\log 1\/\\varepsilon)$ that computes a one-ended spanning forest $\\mathcal{T}$ on $T_\\Delta$.\nMoreover, every vertex $v\\in T_\\Delta$ can compute where it is matched with an uniform local complexity of $O(\\operatorname{poly} \\log 1\/\\varepsilon)$.\nThis follows from the discussion above, as a vertex $v$ can determine where it is matched once it has computed $\\mathcal{T}^{\\leftarrow}(w)$ for each $w$ in its neighborhood.\nHence, $\\Pi^2_{\\text{pm}}$ is in the class $\\mathsf{OLOCAL}(O(\\operatorname{poly} \\log 1\/\\varepsilon))$.\n\nSimilarly, the one-ended spanning forest $\\mathcal{T}\\subseteq \\mathcal{G}$ can be constructed in the class $\\mathsf{MEASURE}$ by \\cite{BrooksMeas}.\nBy the discussion above, this immediately implies that $\\Pi^2_{\\text{pm}}$ is in the class $\\mathsf{MEASURE}$. \n\\end{proof}\n\n\n\n\n\\subsubsection{Decomposition in $\\mathsf{BOREL}$}\nNext, we show that a similar, but weaker, decomposition can be done in a Borel way.\nNamely, we say that a subset of edges of an acyclic infinite graph $G$, denoted by $\\mathcal{T}$, is a \\emph{one or two-ended spanning forest} if every vertex $v\\in G$ is contained in an infinite connected component of $\\mathcal{T}$ and each infinite connected component of $\\mathcal{T}$ has exactly one or two directions to infinity.\nLet $S$ be such a connected component.\nNote that if $S$ has one end, then we can solve $\\Pi^2_{\\text{pm}}$ on $S$ as above.\nIf $S$ has two ends, then $S$ contains a doubly infinite path $P$ with the property that erasing $P$ splits $S\\setminus P$ into finite connected components.\nIn order to solve $\\overline{\\Pi^2_{\\text{pm}}}$ we simply solve $\\Pi^2_{\\text{pm}}$ on the finite connected components of $S\\setminus P$ (with some of the vertices in $S\\setminus P$ being potentially matched with vertices on the path $P$) and declare vertices on $P$ to be line vertices.\n\nThe high-level idea to find a one or two-ended spanning forest is to do the same construction as in the measure case and understand what happens with edges that do not disappear after countably many steps.\nA slight modification in the construction guarantees that what remains are doubly infinite paths.\n\n\\begin{restatable}{theorem}{oneendedforestborel}\n\\label{thm:OneTwoEndBorel}\nLet $\\mathcal{G}$ be a Borel $\\Delta$-regular forest. \nThen there is a Borel one or two-ended spanning forest $\\mathcal{T}\\subseteq \\mathcal{G}$.\n\\end{restatable}\n\nThe proof of \\cref{thm:OneTwoEndBorel} can be found in \\cref{sec:oneortwo}.\nWe remark that the notation used in the proof is close to the notation in the proof that gives a measurable construction of a one-ended spanning forest in \\cite{BrooksMeas}.\nNext, we show more formally how \\cref{thm:OneTwoEndBorel} implies that $\\overline{\\Pi^2_{\\text{pm}}}\\in \\mathsf{BOREL}$.\n\n\\begin{theorem}\nThe problem $\\overline{\\Pi^2_{\\text{pm}}}$ is in the class $\\mathsf{BOREL}$.\n\\end{theorem}\n\\begin{proof}\n\nLet $\\mathcal{T}$ be the one or two-ended spanning forest given by \\cref{thm:OneTwoEndBorel} and $S$ be a connected component of $\\mathcal{T}$.\nIf $S$ is one-ended, then we use the first part of the proof of \\cref{thm:twomatchFFIID} to find a solution of $\\Pi^2_{\\text{pm}}$ on $S$.\nSince deciding that $S$ is one-ended as well as computing an orientation towards infinity can be done in a Borel way, this yields a Borel labeling.\n\nSuppose that $S$ is two-ended.\nThen, there exists a doubly infinite path $P$ in $S$ with the property that the connected components of $S\\setminus P$ are finite.\nMoreover, it is possible to detect $P$ in a Borel way.\nThat is, declaring the vertices on $P$ to be line vertices and using the special symbol in $\\overline{\\Pi^2_{\\text{pm}}}$ for the half-edges on $P$ yields a Borel measurable labeling.\nLet $C\\not=\\emptyset$ be one of the finite components in $S\\setminus P$ and $v\\in C$ be the vertex of distance $1$ from $P$.\nOrient the edges in $C$ towards $v$, this can be done in a Borel way since $v$ is uniquely determined for $C$.\nThen the first part of the proof of \\cref{thm:twomatchFFIID} shows that one can inductively find a solution to $\\Pi^2_{\\text{pm}}$ on $C$ in such a way that all vertices, possibly up to $v$, are matched.\nIf $v$ is matched, then we are done.\nIf $v$ is not matched, then we add the edge that connects $v$ with $P$ to the matching. Note that it is possible that multiple vertices are matched with the same path vertex, but this is not a problem according to the definition of $\\overline{\\Pi^2_{\\text{pm}}}$. \nIt follows that this defines a Borel function on $H(\\mathcal{G})$ that solves $\\overline{\\Pi^2_{\\text{pm}}}$.\n\\end{proof}\n\n\n\n\n\n\\subsubsection{Vizing's Theorem for $\\Delta=3$}\nFinding a measurable or local version of Vizing's Theorem, i.e., proper edge $(\\Delta+1)$-coloring $\\Pi_{\\chi',\\Delta+1}$, was studied recently in \\cite{grebik2020measurable,weilacher2021borel,BernshteynVizing}.\nIt is however not known, even on trees, whether $\\Pi_{\\chi',\\Delta+1}$ is in $\\mathsf{RLOCAL}(O(\\log\\log n))$.\nHere we use the one-ended forest construction to show that $\\Pi_{\\chi',\\Delta+1}$ is in the class $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon)$ for $\\Delta=3$.\n\n\n\n\\iffalse\n\\paragraph{Composition of Uniform Local Algorithms}\nWhen designing algorithms, one often composes several simple algorithms to obtain a more sophisticated algorithm. For example, an algorithm $\\mathcal{A}$ might first invoke algorithm $\\mathcal{A}_1$ and after $\\mathcal{A}_1$ finished, $\\mathcal{A}$ invokes algorithm $\\mathcal{A}_2$ with the input of $\\mathcal{A}_2$ being equal to the output of $\\mathcal{A}_1$. Then, the final output of $\\mathcal{A}$ is simply the output of $\\mathcal{A}_2$. We call such a composition a sequential composition. In the context of the classic deterministic local model, the round complexity of $\\mathcal{A}$ is simply the sum of the round complexities of the two algorithms $\\mathcal{A}_1$ and $\\mathcal{A}_2$ that are used as a subroutine. Here, we prove two such composition results in the context of uniform local algorithms. These general results will help us in \\cref{app:one_ended} to analyze the uniform complexity of an uniform algorithm that computes a so-called one-ended forest in a principled manner.\n\n\n\n\\begin{lemma}[Parallel Composition]\n\\label{lem:parallel_composition}\nLet $A_1$ and $A_2$ be two distributed uniform algorithms with an uniform local complexity of $t_{A_1}(\\varepsilon)$ and $t_{A_2}(\\varepsilon)$, respectively. \nLet $A$ be the parallel composition of $A_1$ and $A_2$. That is, the coding radius of algorithm $\\mathcal{A}$ at a vertex $u$ is the smallest radius such that $A$ knows the local answer of both $A_1$ and $A_2$ at $u$. \nThen, $t_A(\\varepsilon) \\leq \\max(t_{A_1}(\\varepsilon\/2), t_{A_2}(\\varepsilon\/2))$.\n\\end{lemma}\n\\begin{proof}\nThe probability that the coding radius of algorithm $\\mathcal{A}_1$ at vertex $u$ is larger than $\\max(t_{A_1}(\\varepsilon\/2), t_{A_2}(\\varepsilon\/2))$ is at most $\\varepsilon\/2$ and the same holds for the coding radius of algorithm $\\mathcal{A}_2$ at vertex $u$. Hence, by a simple union bound, the probability that the coding radius of algorithm $\\mathcal{A}$ at a vertex $u$ is larger than $\\max(t_{A_1}(\\varepsilon\/2), t_{A_2}(\\varepsilon\/2))$ is at most $\\varepsilon\/2 + \\varepsilon\/2 = \\varepsilon$. \\\\ Hence, $t_A(\\varepsilon) \\leq \\max(t_{A_1}(\\varepsilon\/2), t_{A_2}(\\varepsilon\/2))$, as desired.\n\\end{proof}\n\n\\fi\n\n\n\n\\begin{proposition}\nLet $\\Delta=3$. We have $\\Pi_{\\chi',\\Delta+1} \\in\\mathsf{OLOCAL}(\\operatorname{poly} \\log 1\/\\varepsilon)$.\n\\end{proposition}\n\\begin{proof}[Proof sketch.]\nBy \\cref{thm:one_ended_forest}, we can compute a one-ended forest decomposition $\\mathcal{T}$ with uniform local complexity $O(\\operatorname{poly} \\log 1\/\\varepsilon)$. \nNote that every vertex has at least one edge in $\\mathcal{T}$, hence the edges in $G \\setminus \\mathcal{T}$ form paths. \nThese paths can be $3$-edge colored by using the uniform version of Linial's coloring algorithm. This algorithm has an uniform local complexity of $O(\\log^* 1\/\\varepsilon)$. By \\cref{lem:sequential_composition}, the overall complexity is $O(\\operatorname{poly} \\log (\\Delta^{\\log^* 1\/\\varepsilon} \/ \\varepsilon) + \\log^* 1\/\\varepsilon)) = O(\\operatorname{poly} \\log 1\/\\varepsilon)$. \n\nFinally, we color the edges of $T$. We start from the leaves and color the edges inductively. In particular, whenever we consider a vertex $v$ we color the at most two edges in $T^{\\leftarrow}(v)$ below $v$. \nNote that there always is at least one uncolored edge going from $v$ in the direction of $T$. Hence, we can color the at most two edges below $v$ in $T^{\\leftarrow}(v)$ greedily -- each one neighbors with at most $3$ colored edges at any time. \n\\end{proof}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{Proof of \\cref{thm:one_ended_forest}} \n\\label{sec:formalproof}\n\n\nIn this section we formally prove \\cref{thm:one_ended_forest}. \n\n\n\\oneendedtreetail*\n\nWe follow the construction of the one-ended spanning forest in \\cite{BrooksMeas}.\nThe construction proceeds in rounds. Before each round, some subset of the vertices have already decided which incident edge to grab, and we refer to these vertices as settled vertices. All the remaining vertices are called unsettled. The goal in each round is to make a large fraction of the unsettled vertices settled. To achieve this goal, our construction relies on certain properties that the graph induced by all the unsettled vertices satisfies. \n\nOne important such property is that the graph induced by all the unsettled vertices is expanding. That is, the number of vertices contained in the neighborhood around a given vertex grows exponentially with the radius. The intuitive reason why this is a desirable property is the following. Our algorithm will select a subset of the unsettled vertices and clusters each unsettled vertex to the closest selected vertex. As the graph is expanding and any two vertices in the selected subset are sufficiently far away, there will be a lot of edges leaving a given cluster. For most of these inter-cluster edges, both of the endpoints will become settled. This in turn allows one to give a lower bound on the fraction of vertices that become settled in each cluster.\n\nTo ensure that the graph induced by all the unsettled vertices is expanding, a first condition we impose on the graph induced by all the unsettled vertices is that it has a minimum degree of at least $2$. While this condition is a first step in the right direction, it does not completely suffice to ensure the desired expansion, as an infinite path has a minimum degree of $2$ but does not expand sufficiently. Hence, our algorithm will keep track of a special subset of the unsettled vertices, the so-called hub vertices. Each hub vertex has a degree of $\\Delta$ in the graph induced by all the unsettled vertices. That is, all of the neighbors of a hub vertex are unsettled as well. Moreover, each unsettled vertex has a bounded distance to the closest hub vertex, where the specific upper bound on the distance to the closest hub vertex increases with each round. As we assume $\\Delta > 2$, the conditions stated above suffice to show that the graph induced by all the unsettled vertices expands.\n\nNext, we explain in more detail how a vertex decides which edge to grab. Concretely, if a vertex becomes settled, it grabs an incident edge such that the other endpoint of that edge is strictly closer to the closest unsettled vertex as the vertex itself. For example, if a vertex becomes settled and there is exactly one neighbor that is still unsettled, then the vertex will grab the edge that the vertex shares with its unsettled neighbor. Grabbing edges in that way, one can show two things. First, for a settled vertex $v$, the set $T^{\\leftarrow}(v)$ of vertices behind $v$ does not change once $v$ becomes settled. The intuitive reason for this is that the directed edges point towards the unsettled vertices and therefore the unsettled vertices lie before $v$ and not after $v$. Moreover, one can also show that $T^{\\leftarrow}(v)$ only contains finitely many vertices. The reason for this is that at the moment a vertex $v$ becomes settled, there exists an unsettled vertex that is sufficiently close to $v$, where the exact upper bound on that distance again depends on the specific round in which $v$ becomes settled. \n\nWhat remains to be discussed is at which moment a vertex decides to become settled. As written above, in each round the algorithm considers a subset of the unsettled vertices and each unsettled vertex is clustered to the closest vertex in that subset. This subset only contains hub vertices and it corresponds to an MIS on the graph with the vertex set being equal to the set of hub vertices and where two hub vertices are connected by an edge if they are sufficiently close in the graph induced by all the unsettled vertices. Now, each cluster center decides to connect to exactly $\\Delta$ different neighboring clusters, one in each of its $\\Delta$ subtrees. Remember that as each cluster center is a hub vertex, all of its neighbors are unsettled as well. Now, each unsettled vertex becomes settled except for those that lie on the unique path between two cluster centers such that one of the cluster centers decided to connect to the other one.\n\n\n\n\n\n\n\\paragraph{Construction}\nThe algorithm proceeds in rounds. After each round, a vertex is either settled or unsettled and a settled vertex remains settled in subsequent rounds. Moreover, some unsettled vertices are so-called hub vertices. We denote the set of unsettled, settled and hub vertices after the $i$-th round with $U_i$, $S_i$ and $H_i$, respectively. We set $U_0 = H_0 = V$ prior to the first round and we always set $S_i = V \\setminus U_i$. \\cref{fig:constr1} illustrates the situation after round $i$. Moreover, we denote with $O_i$ a partial orientation of the edges of the infinite $\\Delta$-regular input tree $T$. In the beginning, $O_0$ corresponds to the partial orientation with no oriented edges. Each vertex is incident to at most $1$ outwards oriented edge in $O_i$. If a vertex is incident to an outwards oriented edge in $O_i$, then the other endpoint of the edge will be its parent in the one-ended forest decomposition. For each vertex $v$ in $S_i$, we denote with $T_{v,i}$ the smallest set that satisfies the following conditions. First, $v \\in T_{v,i}$. Moreover, if $u \\in T_{v,i}$ and $\\{w,u\\}$ is an edge that according to $O_i$ is oriented from $w$ to $u$, then $w \\in T_{v,i}$. We later show that $T_{v,i}$ contains exactly those vertices that are contained in the subtree $\\mathcal{T}^\\leftarrow(v)$.\n\n\n\\begin{figure}[H]\n \\centering\n \\includegraphics{fig\/one_ended_construction.pdf}\n \\caption{The picture illustrates the situation after some round $i$ of the construction. The vertices on the paths are all contained in $U_i$, while the vertices corresponding to the big black dots are additionally contained in $H_i$. Note that the length of each path is at least $d_i$. However, the distance between $u$ and $v$ in the original graph can be much smaller.}\n \\label{fig:constr1}\n\\end{figure}\n\n\nAfter the $i$-th round, the construction satisfies the following invariants.\n\n\n\n\\begin{enumerate}\n \n \\item For each $v \\in S_{i-1}$, we have $T_{v,i} = T_{v,i-1}$.\n \n \\item Let $v$ be an arbitrary vertex in $S_i$. Then, $T_{v,i}$ contains finitely many vertices and furthermore $T_{v,i} \\subseteq S_i$.\n \n \\item The minimum degree of the graph $T[U_i]$ is at least $2$ and each vertex in $H_i$ has a degree of $\\Delta$ in $T[U_i]$.\n \n \\item Each vertex in $U_i$ has a distance of at most $\\sum_{j=0}^i d_j$ to the closest vertex in $H_i$ in the graph $T[U_i]$.\n\\end{enumerate}\n\nWe now describe how to compute $U_i, S_i$, $H_i$ and $O_i$. The construction during the $i$-th round is illustrated in $\\cref{fig:constr2}$. Note that we can assume that the invariants stated above are satisfied after the $(i-1)$-th round. In the $i$-th round we have a parameter $d_i$ that we later set to $2^{2^i}$.\n\n\\begin{enumerate}\n \\item $H_i$ is a subset of $H_{i-1}$ that satisfies the following property. No two vertices in $H_i$ have a distance of at most $d_i$ in $T[U_{i-1}]$. Moreover, each vertex in $H_{i-1} \\setminus H_i$ has a distance of at most $d_i$ to the closest vertex in $H_i$ in the graph $T[U_{i-1}]$. Note that we can compute $H_i$ by computing an MIS in the graph with vertex set $H_{i-1}$ and where two vertices are connected iff they have a distance of at most $d_i$ in the graph $T[U_{i-1}]$. Hence, we can use Ghaffari's MIS algorithm \\cite{ghaffari2016MIS} to compute the set $H_i$.\n \\item Next, we describe how to compute $U_i$. \n We assign each vertex $u \\in U_{i-1}$ to the closest vertex in the graph $T[U_{i-1}]$ that is contained in $H_i$, with ties being broken arbitrarily. We note that there exists a node in $H_i$ with a distance of at most $(\\sum_{j=0}^{i-1} d_j) + d_i$ to $u$. To see why, note that Invariant $(4)$ from round $i - 1$ implies that there exists a vertex $w$ in $H_{i-1}$ with a distance of at most $\\sum_{j=0}^{i-1} d_j$ to $u$. We are done if $w$ is also contained in $H_i$. If not, then it follows from the way we compute $H_i$ that there exists a vertex in $H_i$ with a distance of at most $d_i$ to $w$, but then the triangle inequality implies that the distance from $u$ to that vertex is at most $(\\sum_{j=0}^{i-1} d_j) + d_i$, as desired. For each vertex $v \\in H_i$, we denote with $C(v)$ the set of all vertices in $U_{i-1}$ that got assigned to $v$. Now, let $E_v$ denote the set of edges that have exactly one endpoint in $C(v)$. We can partition $E_v$ into $E_{v,1} \\sqcup E_{v,2} \\sqcup \\ldots \\sqcup E_{v,\\Delta}$, where for $\\ell \\in [\\Delta]$, $E_{v,\\ell}$ contains all the edges in $E_v$ that are contained in the $\\ell$-th subtree of $v$. \n Invariants $(3)$ and $(4)$ after the $(i-1)$-th round imply that $E_{v,\\ell} \\neq \\emptyset$. Moreover, $E_{v,\\ell}$ contains only finitely many edges, as we have shown above that the cluster radius is upper bounded by $\\sum_{j=0}^{i} d_j$.\n \n Now, for $\\ell \\in [\\Delta]$, we choose an edge $e_\\ell$ uniformly at random from the set $E_{v,\\ell}$. Let $u_\\ell \\neq v$ be the unique vertex in $H_i$ such that one endpoint of $e_\\ell$ is contained in $C(u_\\ell)$. We denote with $P_{v,\\ell}$ the set of vertices that are contained in the unique path between $v$ and $u_\\ell$. Finally, we set $U_i = \\bigcup_{v \\in H_i, \\ell \\in [\\Delta]} P_{v,\\ell}$.\n \\item It remains to describe how to compute the partial orientation $O_i$. All edges that are oriented in $O_{i-1}$ will be oriented in the same direction in $O_i$. Additionally, we orient for each vertex $u$ that got settled in the $i$-th round, i.e., $u \\in S_i \\cap U_{i-1}$, exactly one incident edge away from $u$. Let $w$ be the closest vertex to $u$ in $T[U_{i-1}]$ that is contained in $U_i$, with ties being broken arbitrarily. We note that it follows from the discussions above that the distance between $u$ and $w$ in the graph $T[U_{i-1}]$ is at most $\\sum_{j=0}^{i} d_j$. We now orient the edge incident to $u$ that is on the unique path between $u$ and $w$ outwards. We note that this orientation is well-defined in the sense that we only orient edges that were not oriented before and that we don't have an edge such that both endpoints of that edge want to orient the edge outwards.\n\\end{enumerate}\n\n\n\\begin{figure}\n \\centering\n \\includegraphics{fig\/one_ended_construction2.pdf}\n \\caption{The picture illustrates the situation during some round $i$ of the construction. The vertices on all paths are contained in $U_{i-1}$ and the vertices corresponding to the big black dots are contained in $H_{i-1}$. The vertices $v$, $u_1$, $u_2$ and $u_3$ are the only vertices in the picture that are contained in $H_i$. All the vertices on the red paths remain unsettled during round $i$ while all the vertices on the black paths become settled. Each vertex that becomes settled during the $i$-th round orients one incident edge outwards, illustrated by the green edges. Each green edge is oriented towards the closest vertex that is still unsettled after the $i$-th round. We note that some of the edges in the picture are oriented away from $v$.}\n \\label{fig:constr2}\n\\end{figure}\n\n\nWe now prove by induction that our procedure satisfies all the invariants.\nPrior to the first round, all invariants are trivially satisfied. Hence, it remains to show that the invariants hold after the $i$-th round, given that the invariants hold after the $(i-1)$-th round. \n\n\\begin{enumerate}\n\n \\item Let $v \\in S_{i-1}$ be arbitrary. As $O_i$ is an extension of $O_{i-1}$, it follows from the definition of $T_{v,i-1}$ and $T_{v,i}$ that $T_{v,i-1} \\subseteq T_{v,i}$. Now assume that $T_{v,i} \\not \\subseteq T_{v,i-1}$. This would imply the existence of an edge $\\{u,w\\}$ such that $u \\in T_{v,i-1}$ and the edge was oriented during the $i$-th round from $w$ to $u$. As $u \\in T_{v,i-1}$, Invariant $(2)$ from the $(i-1)$-th round implies $u \\in S_{i-1}$. This is a contradiction as during the $i$-th round only edges with both endpoints in $U_{i-1}$ can be oriented. Therefore it holds that $T_{v,i} = T_{v,i-1}$, as desired.\n\n \\item Let $v \\in S_i$ be arbitrary. If it also holds that $v \\in S_{i-1}$, then the invariant from round $i-1$ implies that $T_{v,i-1} = T_{v,i}$ contains finitely many vertices and $T_{v,i-1} \\subseteq S_{i-1} \\subseteq S_i$. Thus, it suffices to consider the case that $v \\in S_i \\cap U_{i-1}$. Note that $T_{v,i} \\cap U_i = \\emptyset$. Otherwise there would exist an edge that is oriented away from a vertex in $U_i$ according to $O_i$, but this cannot happen according to the algorithm description. Hence, $T_{v,i} \\subseteq S_i$. Now, for the sake of contradiction, assume that $T_{v,i}$ contains infinitely many vertices. From the definition of $T_{v,i}$, this implies that there exists a sequence of vertices $(v_k)_{k \\geq 1}$ with $v = v_1$ such that for each $k \\geq 1$, $\\{v_{k+1},v_k\\}$ is an edge in $T$ that is oriented from $v_{k+1}$ to $v_k$ according to $O_i$. For each $k \\geq 1$ we have $v_k \\in S_i$. We furthermore know that $v_k \\in U_{i-1}$, as otherwise $T_{v_k,i} = T_{v_k,i-1}$ would contain infinitely many vertices, a contradiction. From the way we orient the edges, a simple induction proof implies that the closest vertex of $v_k$ in the graph $T[U_{i-1}]$ that is contained in $U_i$ has a distance of at least $k$. However, we previously discussed that the distance between $v_k$ and the closest vertex in $U_i$ in the graph $T[U_{i-1}]$ is at most $\\sum_{j=0}^{i} d_j$. This is a contradiction. Hence, $T_{v,i}$ contains finitely many vertices.\n \n \\item It follows directly from the description of how to compute $U_i$ and $H_i$ that the minimum degree of the graph $T[U_i]$ is at least $2$ and that each vertex in $H_i$ has a degree of $\\Delta$ in $T[U_i]$.\n \n \\item Let $u \\in U_i$ be arbitrary. We need to show that $u$ has a distance of at most $\\sum_{j=0}^i d_j$ to the closest vertex in $H_i$ in the graph $T[U_i]$. Let $v$ be the vertex with $u \\in C(v)$. We know that there is no vertex in $H_i$ that is closer to $u$ in $T[U_{i-1}]$ than $v$. Hence, the distance between $u$ and $v$ is at most $\\sum_{j=0}^i d_j$ in the graph $T[U_{i-1}]$. Moreover, from the description of how we compute $U_i$, it follows that all the vertices on the unique path between $u$ and $v$ are contained in $U_i$. Hence, each vertex in $U_i$ has a distance of at most $\\sum_{j=0}^i d_j$ to the closest vertex in $H_i$ in the graph $T[U_i]$, as desired.\n \\end{enumerate}\n\nWe derive an upper bound on the coding radius of the algorithm in three steps. First, for each $i \\in \\mathbb{N}$ and $\\varepsilon > 0$, we derive an upper bound on the coding radius for computing with probability at least $1 - \\varepsilon$ for a given vertex in which of the sets $S_i, U_i$ and $H_i$ it is contained in and for each incident edge whether and how it is oriented in $O_i$, given that each vertex receives as additional input in which of the sets $S_{i-1}, U_{i-1}$ and $H_{i-1}$ it is contained in and for each incident edge whether and how it is oriented in $O_{i-1}$. Given this upper bound on the coding radius, we can use the sequential composition lemma (\\cref{lem:sequential_composition}) and a simple induction proof to give an upper bound on the coding radius for computing with probability at least $1 - \\varepsilon$ for a given vertex in which of the sets $S_i, U_i$ and $H_i$ it is contained in and for each incident edge whether and how it is oriented in $O_i$, this time without providing any additional input. \n\nSecond, we analyze after how many rounds a given vertex is settled with probability at least $1- \\varepsilon$ for a given $\\varepsilon > 0$. \n\nFinally, we combine these two upper bounds to prove \\cref{thm:one_ended_forest}.\n \n \n\\begin{lemma}\n\\label{lem:tail1}\nFor each $i \\in \\mathbb{N}$ and $\\varepsilon \\in (0,0.01]$, let $g_i(\\varepsilon)$ denote the smallest coding radius such that with probability at least $1 - \\varepsilon$ one knows for a given vertex in which of the sets $S_i, U_i$ and $H_i$ it is contained in and for each incident edge whether and how it is oriented in $O_i$, given that each vertex receives as additional input in which of the sets $S_{i-1}, U_{i-1}$ and $H_{i-1}$ it is contained in and for each incident edge whether and how it is oriented in $O_{i-1}$. It holds that $g_i(\\varepsilon) = O(d^2_i \\log (\\Delta\/\\varepsilon))$.\n\\end{lemma}\n\n\\begin{proof}\n When running Ghaffari's uniform MIS algorithm on some graph $G'$ with maximum degree $\\Delta'$, each vertex knows with probability at least $1-\\varepsilon$ whether it is contained in the MIS or not after $O(\\log (\\Delta') + \\log (1\/\\varepsilon))$ communication rounds in $G'$ (\\cite{ghaffari2016MIS}, Theorem 1.1). For the MIS computed during the $i$-th round, $\\Delta' =\\Delta^{O(d_i)}$ and each communication round in $G'$ can be simulated with $O(d_i)$ communication rounds in the tree $T$. Hence, to know for a given vertex with probability at least $1-\\varepsilon$ whether it is contained in the MIS or not, it suffices consider the $O(d_i^2\\log(\\Delta) + d_i\\log(1\/\\varepsilon))$-hop neighborhood around that vertex. Now, let $u$ be an arbitrary vertex. In order to compute in which of the sets $S_i, U_i$ and $H_i$ $u$ is contained in and for each incident edge of $u$ whether and how it is oriented in $O_i$, we not only need to know whether $u$ is in the MIS or not. However, it suffices if we know for all the vertices in the $O(d_i)$-hop neighborhood of $u$ whether they are contained in the MIS or not (on top of knowing for each vertex in the neighborhood in which of the sets $S_{i-1}, U_{i-1}$ and $H_{i-1}$ it is contained in and for each incident edge whether and how it is oriented in $O_{i-1}$). Hence, by a simple union bound over the at most $\\Delta^{O(d_i)}$ vertices in the $O(d_i)$-hop neighborhood around $u$, we obtain $g_i(\\varepsilon) = O(d_i) + O(d_i^2\\log(\\Delta) + d_i\\log(\\Delta^{O(d_i)}\/\\varepsilon)) = O(d_i^2 \\log(\\Delta\/\\varepsilon))$.\n\\end{proof}\n\n\\begin{lemma}\n\\label{lem:tail2}\nFor each $i \\in \\mathbb{N}$ and $\\varepsilon \\in (0,0.01]$, let $h_i(\\varepsilon)$ denote the smallest coding radius such that with probability at least $1 - \\varepsilon$ one knows for a given vertex in which of the sets $S_i, U_i$ and $H_i$ it is contained in and for each incident edge whether and how it is oriented in $O_i$. Then, there exists a constant $c$ independent of $\\Delta$ such that $h_i(\\varepsilon) \\leq 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta\/\\varepsilon)$.\n\\end{lemma}\n\n\\begin{proof}\nBy prove the statement by induction on $i$. For a large enough constant $c$, it holds that $h_1(\\varepsilon) \\leq 2^{2^{3}} \\cdot(c\\log(\\Delta)) \\log(\\Delta\/\\varepsilon)$. Now, consider some arbitrary $i$ and assume that $h_i(\\varepsilon) \\leq 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta\/\\varepsilon)$ for some large enough constant $c$. We show that this implies $h_{i+1}(\\varepsilon) \\leq 2^{2^{(i+1)+2}} \\cdot(c\\log(\\Delta))^{i+1} \\log(\\Delta\/\\varepsilon)$.\nBy the sequential composition lemma (\\cref{lem:sequential_composition}) and assuming that $c$ is large enough, we have\n\n\\begin{align*}\n h_{i+1}(\\varepsilon) &\\leq h_{i}((\\varepsilon\/2)\/\\Delta^{g_{i+1}(\\varepsilon\/2) + 1}) + g_{i+1}(\\varepsilon\/2) \\\\\n &\\leq 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta \\cdot \\Delta^{g_{i+1}(\\varepsilon\/2) + 1}\/(\\varepsilon\/2)) + g_{i+1}(\\varepsilon\/2) \\\\\n &\\leq 2 \\cdot 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta \\cdot \\Delta^{g_{i+1}(\\varepsilon\/2) + 1}\/(\\varepsilon\/2)) \\\\\n &\\leq 2 \\cdot 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta^{(c\/10) \\cdot 2^{2^{i+2}}\\log(\\Delta\/\\varepsilon)}\/(\\varepsilon\/2)) \\\\\n &\\leq 4 \\cdot 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta^{(c\/10) \\cdot 2^{2^{i+2}}\\log(\\Delta\/\\varepsilon)}) \\\\\n &\\leq (4\/10) \\cdot 2^{2^{i+2}} \\cdot 2^{2^{i+2}} (c \\log(\\Delta))^{i+1} \\log(\\Delta\/\\varepsilon) \\\\\n &\\leq 2^{2^{(i+1)+2}} \\cdot(c\\log(\\Delta))^{i+1} \\log(\\Delta\/\\varepsilon),\n\\end{align*}\n\nas desired.\n\n\\end{proof}\n\n\n \n \n\n\\begin{lemma}\\label{lm:tail0}\nLet $u$ be an arbitrary vertex. For each $\\varepsilon \\in (0,0.01]$, let $f(\\varepsilon)$ denote the smallest $i \\in \\mathbb{N}$ such that $u$ is settled after the $i$-th round with probability at least $1 - \\varepsilon$. There exists a fixed $c \\in \\mathbb{R}$ independent of $\\Delta$ such that $f(\\varepsilon) \\leq \\lceil 1 + \\log \\log \\frac{1}{c} \\log_\\Delta(1\/\\varepsilon)\\rceil$.\n\\end{lemma}\n\n\\begin{proof}\n\nLet $i \\in \\mathbb{N}$ be arbitrary. We show that a given vertex is settled after the $i$-th round with probability at least $1 - O(1\/\\Delta^{\\Omega(d_i\/d_{i-1})})$. For the sake of analysis, we run the algorithm on a finite $\\Delta$-regular high-girth graph instead of an infinite $\\Delta$-regular tree. For now, we additionally assume that no vertex realizes that we don't run the algorithm on an infinite $\\Delta$-regular tree. That is, we assume that for each vertex the coding radius to compute all its local information after the $i$-th round is much smaller than the girth of the graph. \n\nWith this assumption, we give a deterministic upper bound on the fraction of vertices that are not settled after the $i$-th round. On the one hand, the number of vertices that are not settled after the $i$-th round can be upper bounded by $|H_i| \\cdot \\Delta \\cdot O(d_i)$. On the other hand, we will show that the fraction of vertices that are contained in $H_i$ is upper bounded by $1\/\\Delta^{\\Omega(d_i\/d_{i-1})}$. We do this by showing that $C(v)$ contains $\\Delta^{\\Omega(d_i\/d_{i-1})}$ many vertices for a given $v \\in H_i$. Combining these two bounds directly implies that the fraction of unsettled vertices is smaller than \n\n\\[\\frac{\\Delta \\cdot O(d_i)}{\\Delta^{\\Omega(d_i\/d_{i-1})}} = 1\/\\Delta^{\\Omega(d_i\/d_{i-1})}.\\]\n\n\nLet $v \\in H_i$ be arbitrary. We show that $|C(v)| = \\Delta^{\\Omega(d_i\/d_{i-1})}$. Using Invariants (3) and (4) together with a simple induction argument, one can show that there are at least $(\\Delta-1)^{\\lfloor D \/ ((2\\sum_{j=0}^{i-1}d_{j}) + 1)\\rfloor}$ vertices contained in $H_{i-1}$ and whose distance to $v$ is at most $D$ in the graph induced by the vertices in $U_{i-1}$. Furthermore, from the way we defined the clustering $C(v)$ and the fact that two vertices in $H_i$ have a distance of at least $d_i$ in the graph $T[U_{i-1}]$, it follows that all vertices in $U_{i-1}$ having a distance of at most $d_i\/2 -1$ to $v$ are contained in the cluster $C(v)$. Hence, the total number of vertices in $C(v)$ is at least $(\\Delta-1)^{\\lfloor (d_i\/2 - 1) \/ ((2\\sum_{j=0}^{i-1}d_{j}) + 1)\\rfloor} = \\Delta^{\\Omega(d_i\/d_{i-1})}$, as promised.\n\nNow we remove the assumption that no vertex realizes that we don't run the algorithm on an infinite $\\Delta$-regular tree. If a vertex does realize that we don't run the algorithm on an infinite $\\Delta$-regular tree, then we consider the vertex as being unsettled after the $i$-th round. By considering graphs with increasing girth, we can make the expected fraction of vertices that realize that we are not on an infinite $\\Delta$-regular tree arbitrarily small. Combining this observation with the previous discussion, this implies that the expected fraction of vertices that are not settled after the $i$-th round is at most $1\/\\Delta^{\\Omega(d_i\/d_{i-1})}$. By symmetry, each vertex has the same probability of being settled after the $i$-th round. Hence, the probability that a given vertex is settled after the $i$-th round when run on a graph with sufficiently large girth is $1 - 1\/\\Delta^{\\Omega(d_i\/d_{i-1})}$, and the same holds for each vertex when we run the algorithm on an infinite $\\Delta$-regular tree. Thus, there exists a constant $c$ such that the probability that a given vertex is unsettled after the $i$-th round is at most $1\/\\Delta^{c \\cdot d_i\/d_{i-1}} = 1\/\\Delta^{c \\cdot 2^{2^{i-1}}}$. Setting $i = \\lceil 1 + \\log \\log \\frac{1}{c} \\log_\\Delta(1\/\\varepsilon)\\rceil$ finishes the proof.\n\\end{proof}\n\n\n\n\n\nWe are now finally ready to finish the proof of \\cref{thm:one_ended_forest}. Let $u$ be an arbitrary vertex and $\\varepsilon \\in (0,0.01]$. We need to compute an upper bound on the coding radius that is necessary for $u$ to know all the vertices in $\\mathcal{T}^{\\leftarrow}(u)$ with probability at least $1-\\varepsilon$. To compute such an upper bound for the required coding radius it suffices to find an $i^*$ and an $R^*$ such that the following holds. First, $u$ is settled after $i^*$ rounds with probability at least $1 - \\varepsilon\/2$. Second, $u$ knows for each edge in its $O(d_{i^*})$-hop neighborhood whether it is oriented in the partial orientation $O_{i^*}$, and if yes, in which direction, by only considering its $R^*$-hop neighborhood with probability at least $1 - \\varepsilon\/2$. By a union bound, both of these events occur with probability at least $1-\\varepsilon$. Moreover, if both events occur $u$ knows all the vertices in the set $\\mathcal{T}^{\\leftarrow}(u)$. The reason is as follows. If $u$ is settled after round $i^*$, then it follows from the previous analysis that only vertices in the $O(d_{i^*})$-hop neighborhood of $u$ can be contained in $\\mathcal{T}^{\\leftarrow}(u)$. Moreover, for each vertex in its $O(d_{i^*})$-hop neighborhood, $u$ can determine if the vertex is contained in $\\mathcal{T}^{\\leftarrow}(u)$ if it knows all the edge orientations on the unique path between itself and that vertex after the $i^*$-th round. Hence, it remains to find concrete values for $i^*$ and $R^*$. According to \\cref{lm:tail0}, we can choose $i^* = \\lceil 1 + \\log \\log \\frac{1}{c} \\log_\\Delta (2\/\\varepsilon)\\rceil$ for some large enough constant $c$. Moreover, it follows from a union bound that all vertices in the $O(d_{i^*})$-hop neighborhood around $u$ know with probability at least $1-\\varepsilon\/2$ the orientation of all its incident edges according to $O_{i^*}$ by only considering their $h_{i^*}((\\varepsilon\/2)\/\\Delta^{O(d_{i^*})})$-hop neighborhood. Hence, we can set \n\n\\begin{align*}\nR^* &= O(d_{i^*}) + h_{i^*}((\\varepsilon\/2)\/\\Delta^{O(d_{i^*})}) \\\\\n &\\leq O(d_{i^*}) + 2^{2^{i^* + 2}} \\cdot (c \\log(\\Delta))^{i^*}\\log(\\Delta^{O(d_{i^*})}\/\\varepsilon) \\\\\n &= O(2^{2^{i^*}}\\cdot 2^{2^{i^* + 2}} (c \\log(\\Delta))^{i^*}\\log(\\Delta\/\\varepsilon)) \\\\\n &= O(2^{2^{i^* + 3}} (c \\log(\\Delta))^{i^*}\\log(\\Delta\/\\varepsilon)) \\\\\n &= O(2^{2^{\\log \\log \\frac{1}{c} \\log (2\/\\varepsilon) + 5}} (c \\log(\\Delta))^{\\log \\log \\frac{1}{c} \\log (2\/\\varepsilon) + 2}\\log(\\Delta\/\\varepsilon)) \\\\\n &= \\log(1\/\\varepsilon)^{32} \\cdot \\log(\\Delta\/\\varepsilon) \\cdot \\log\\log(1\/\\varepsilon)^{O(\\log\\log(\\Delta))}.\n\\end{align*}\n\nAs $\\Delta = O(1)$, it therefore holds that\n\n\\[P(R(v)) = \\operatorname{poly}(\\log(1\/\\varepsilon)) \\geq 1 - \\varepsilon, \\]\n\nas desired.\n\n\n\n\n\n\n\\section{$\\mathsf{BAIRE} = \\mathsf{LOCAL}(O(\\log n))$}\n\\label{sec:baire}\n\nIn this section, we show that on $\\Delta$-regular trees the classes $\\mathsf{BAIRE}$ and $\\mathsf{LOCAL}(O(\\log(n)))$ are the same.\nAt first glance, this result looks rather counter-intuitive.\nThis is because in finite $\\Delta$-regular trees every vertex can see a leaf of distance $O(\\log(n))$, while there are no leaves at all in an infinite $\\Delta$-regular tree.\nHowever, there is an intuitive reasons why these classes are the same: in both setups there is a technique to decompose an input graph into a hierarchy of subsets. Furthermore, the existence of a solution that is defined inductively with respect to these decompositions can be characterized by the same combinatorial condition of Bernshteyn \\cite{Bernshteyn_work_in_progress}.\nWe start with a high-level overview of the decomposition techniques used in both contexts.\n\n\n\\paragraph{{\\sf Rake}\\ and {\\sf Compress}}\nThe hierarchical decomposition in the context of distributed computing is based on a variant of a decomposition algorithm of\nMiller and Reif~\\cite{MillerR89}.\nTheir original decomposition algorithm works as follows.\nStart with a tree $T$, and repeatedly apply the following two operations alternately: {\\sf Rake}\\ (remove all degree-1 vertices) and {\\sf Compress}\\ (remove all degree-2 vertices).\nThen $O(\\log n)$ iterations suffice to remove all vertices in $T$~\\cite{MillerR89}. To view it another way, this produces a decomposition of the vertex set $V$ into $2L - 1$ layers \\[V = \\VR{1} \\cup \\VC{1} \\cup \\VR{2} \\cup \\VC{2} \\cup \\VR{3} \\cup \\VC{3} \\cup \\cdots \\cup \\VR{L},\\] \nwith $L= O(\\log n)$, where $\\VR{i}$ is the set of vertices removed during the $i$-th {\\sf Rake}\\ operation and $\\VC{i}$ is the set of vertices removed during the $i$-th {\\sf Compress}\\ operation. We will use a variant~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} of this decomposition in the proof of \\cref{thm:ellfull_to_logn}.\n\n\n\n\nVariants of this decomposition turned out to be useful in designing $\\mathsf{LOCAL}$ algorithms~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup,ChangHLPU20,chang2020n1k_speedups}. \nIn our context, we assume that the given LCL satisfies a certain combinatorial condition and then find a solution inductively, in the reversed order of the construction of the decomposition.\nNamely, in the {\\sf Rake}\\ step we want to be able to existentially extend the inductive partial solution to all relative degree $1$-vertices (each $v \\in \\VR{i}$ has degree at most 1 in the subgraph induced by $\\VR{i} \\cup \\cdots \\cup \\VR{L}$) and in the {\\sf Compress}\\ step we want to extend the inductive partial solution to paths with endpoints labeled from the induction (the vertices in $\\VC{i}$ form degree-2 paths in the subgraph induced by $\\VC{i} \\cup \\cdots \\cup \\VR{L}$).\n\n\n\\paragraph{$\\TOAST$}\nFinding a hierarchical decomposition in the context of descriptive combinatorics is tightly connected with the notion of \\emph{Borel hyperfiniteness}.\nUnderstanding what Borel graphs are Borel hyperfinite is a major theme in descriptive set theory \\cite{doughertyjacksonkechris,gaojackson,conley2020borel}.\nIt is known that grids, and generally polynomial growth graphs are hyperfinite, while, e.g., acyclic graphs are not in general hyperfinite \\cite{jackson2002countable}.\nA strengthening of hyperfiniteness that is of interest to us is called a \\emph{toast} \\cite{gao2015forcing,conleymillerbound}.\nA $q$-toast, where $q\\in \\mathbb{N}$, of a graph $G$ is a collection $\\mathcal{D}$ of fintie subsets of $G$ with the property that (i) every pair of vertices is covered by an element of $\\mathcal{D}$ and (ii) the boundaries of every $D\\not=E\\in \\mathcal{D}$ are at least $q$ apart.\nThe idea to use a toast structure to solve LCLs appears in \\cite{conleymillerbound} and has many applications since then \\cite{gao2015forcing,marksunger}.\nThis approach has been formalized in \\cite{grebik_rozhon2021toasts_and_tails}, where the authors introduce $\\TOAST$ algorithms.\nRoughly speaking, an LCL $\\Pi$ admits a $\\TOAST$ algorithm if there is $q\\in \\mathbb{N}$ and a partial extending function (the function is given a finite subset of a tree that is partially colored and outputs an extension of this coloring on the whole finite subset) that has the property that whenever it is applied inductively to a $q$-toast, then it produces a $\\Pi$-coloring.\nAn advantage of this approach is that once we know that a given Borel graph admits, e.g., a Borel toast structure and a given LCL $\\Pi$ admits a $\\TOAST$ algorithm, then we may conclude that $\\Pi$ is in the class $\\mathsf{BOREL}$.\nSimilarly for $\\mathsf{MEASURE}$, $\\mathsf{BAIRE}$ or $\\mathsf{OLOCAL}$, we refer the reader to \\cite{grebik_rozhon2021toasts_and_tails} for more details and results concerning grids.\n\nIn the case of trees there is no way of constructing a Borel toast in general, however, it is a result of Hjorth and Kechris \\cite{hjorth1996borel} that every Borel graph is hyperfinite on a comeager set for every compatible Polish topology.\nA direct consequence of \\cite[Lemma~3.1]{marks2016baire} together with a standard construction of toast via Voronoi cells gives the following strengthening to toast.\nWe include a sketch of the proof for completeness.\n\n\\begin{restatable}{proposition}{BaireToast}\n\\label{pr:BaireToast}\nLet $\\mathcal{G}$ be a Borel graph on a Polish space $(X,\\tau)$ with degree bounded by $\\Delta\\in \\mathbb{N}$.\nThen for every $q>0$ there is a Borel $\\mathcal{G}$-invariant $\\tau$-comeager set $C$ on which $\\mathcal{G}$ admits a Borel $q$-toast.\n\\end{restatable}\n\\begin{proof}[Proof sketch]\nLet $\\{A_n\\}_{n\\in \\mathbb{N}}$ be a sequence of MIS with parameter $f(n)$ as in \\cite[Lemma~3.1]{marks2016baire} for a sufficiently fast growing function $f(n)$, e.g., $f(n)=(2q)^{n^2}$.\nThen $C=\\bigcup_{n\\in \\mathbb{N}}A_n$ is a Borel $\\tau$-comeager set that is $\\mathcal{G}$-invariant.\nWe produce a toast structure in a standard way, e.g., see~\\cite[Appendix~A]{grebik_rozhon2021toasts_and_tails}.\n\nLet $B_n(x)$ denote the ball of radius $f(n)\/3$ and $R_n(x)$ the ball of radius $f(n)\/4$ around $x\\in A_n$.\nIteratively, define cells $\\mathcal{D}_n=\\{C_{n}(x)\\}_{x\\in A_n}$ as follows.\nSet $C_1(x)=R_1(x)$.\nSuppose that $\\mathcal{D}_n$ has been defined and set\n\\begin{itemize}\n\t\\item $H^{n+1}(x,n+1):=R_{n+1}(x)$ for every $x\\in A_{n+1}$,\n\t\\item if $1\\le i\\le n$ and $\\left\\{H^{n+1}(x,i+1)\\right\\}_{x\\in A_{k+1}}$ has been defined, then we put\n\t$$H^{n+1}(x,i)=\\bigcup \\left\\{B_i(y):H^{n+1}(x,i+1)\\cap B_i(y)\\not=\\emptyset\\right\\}$$\n\tfor every $x\\in A_{n+1}$,\n\t\\item set $C_{n+1}(x):=H^{n+1}(x,1)$ for every $x\\in A_{n+1}$, this defines $\\mathcal{D}_{n+1}$.\n\\end{itemize}\nThe fact that $\\mathcal{D}=\\bigcup_{n\\in\\mathbb{N}} \\mathcal{D}_n$ is a $q$-toast on $C$ follows from the fact that $C=\\bigcup_{n\\in \\mathbb{N}}A_n$ together with the fact that the boundaries are $q$ separated which can be shown as in \\cite[Appendix~A]{grebik_rozhon2021toasts_and_tails}.\n\\end{proof}\n\n\n\nTherefore to understand LCLs in the class $\\mathsf{BAIRE}$ we need understand what LCLs on trees admit $\\TOAST$ algorithm.\nIt turns out that these notions are equivalent, again by using the combinatorial characterization of Bernshteyn \\cite{Bernshteyn_work_in_progress} that we now discuss.\n\n\n\\paragraph{Combinatorial Condition -- $\\ell$-full set}\nIn both decompositions, described above, we need to extend a partial coloring along paths that have their endpoints colored from the inductive step.\nThe precise formulation of the combinatorial condition that captures this demand was extracted by Bernshteyn \\cite{Bernshteyn_work_in_progress}.\nHe proved that it characterizes the class $\\mathsf{BAIRE}$ for Cayley graphs of virtually free groups.\nNote that this class contains, e.g., $\\Delta$-regular trees with a proper edge $\\Delta$-coloring.\n\n\\begin{definition}[Combinatorial condition -- an $\\ell$-full set]\n\\label{def:ellfull}\nLet $\\Pi = (\\Sigma, \\mathcal{V}, \\mathcal{E})$ be an LCL and $\\ell\\geq 2$.\nA set $\\mathcal{V}' \\subseteq \\mathcal{V}$ is \\emph{$\\ell$-full} whenever the following is satisfied.\nTake a path with at least $\\ell$ vertices, and add half-edges to it so that each vertex has degree $\\Delta$.\nTake any $c_1, c_2 \\in \\mathcal{V}'$ and label arbitrarily the half-edges around the endpoints with $c_1$ and $c_2$, respectively.\nThen there is a way to label the half-edges around the remaining $\\ell-2$ vertices with configurations from $\\mathcal{V}'$ such that all the $\\ell-1$ edges on the path have valid edge configuration on them. \n\\end{definition}\n\nNow we are ready to formulate the result that combines Bernshteyn's result \\cite{Bernshteyn_work_in_progress} (equivalence between (1.) and (2.), and the moreover part) with the main results of this section.\nThis also shows the remaining implications in \\cref{fig:big_picture_trees}.\n\n\\begin{theorem}\\label{thm:MainBaireLog}\nLet $\\Pi$ be an LCL on regular trees. \nThen the following are equivalent:\n\\begin{enumerate}\n \\item $\\Pi\\in \\mathsf{BAIRE}$,\n \\item $\\Pi$ admits a $\\TOAST$ algorithm,\n \\item $\\Pi$ admits an $\\ell$-full set,\n \\item $\\Pi\\in \\mathsf{LOCAL}(O(\\log(n)))$.\n\\end{enumerate}\nMoreover, any of the equivalent conditions is necessary for $\\Pi\\in \\mathsf{fiid}$.\n\\end{theorem}\n\nNext we discuss the proof of \\cref{thm:MainBaireLog}.\nWe refer the reader to Bernshteyn's paper \\cite{Bernshteyn_work_in_progress} for full proofs in the case of $\\mathsf{BAIRE}$ and $\\mathsf{fiid}$, here we only sketch the argument for completeness.\nWe also note that instead of using the toast construction, he used a path decomposition of acyclic graphs of Conley, Marks and Unger \\cite{conley2020measurable}.\n\n\n\\subsection{Sufficiency}\n\nWe start by showing that the combinatorial condition is sufficient for $\\mathsf{BAIRE}$ and $\\mathsf{LOCAL}(O(\\log(n)))$.\nNamely, it follows from the next results together with \\cref{pr:BaireToast} that (2.) implies all the other conditions in \\cref{thm:MainBaireLog}.\nAs discussed above the main idea is to color inductively along the decompositions.\n\n\\begin{proposition}\\label{pr:toastable}\nLet $\\Pi=(\\Sigma, \\mathcal{V}, \\mathcal{E})$ be an LCL that admits $\\ell$-full set $\\mathcal{V}'\\subseteq \\mathcal{V}$ for some $\\ell>0$.\nThen $\\Pi$ admits a $\\TOAST$ algorithm that produces a $\\Pi$-coloring for every $(2\\ell+2)$-toast $\\mathcal{D}$.\n\\end{proposition}\n\\begin{proof}[Proof sketch]\nOur aim is to build a partial extending function.\nSet $q:=2\\ell+2$.\nLet $E$ be a piece in a $q$-toast $\\mathcal{D}$ and suppose that $D_1,\\dots,D_k\\in \\mathcal{D}$ are subsets of $E$ such that the boundaries are separated.\nSuppose, moreover, that we have defined inductively a coloring of half-edges of vertices in $D=\\bigcup D_i$ using only vertex configurations from $\\mathcal{V}'$ such that every edge configuration $\\mathcal{E}$ is satisfied for every edge in $D$.\n\nWe handle each connected component of $E\\setminus D$ separately.\nLet $A$ be one of them.\nLet $u\\in A$ be a boundary vertex of $E$.\nSuch an vertex exists since every vertex in $E$ has degree $\\Delta$.\nThe distance of $u$ and any $D_i$ is at least $2\\ell+2$ for every $i\\in [k]$.\nWe orient all the edges from $A$ towards $u$.\nMoreover if $v_i\\in A$ is a boundary vertex of some $D_i$ we assign to $v_i$ a path $V_i$ of length $\\ell$ towards $u$.\nNote that $V_i$ and $V_j$ have distance at least $1$, in particular, are disjoint for $i\\not=j\\in [k]$ .\nNow, until you encounter some path $V_i$, color any in manner half-edges of vertices in $A$ inductively starting at $u$ in such a way that edge configurations $\\mathcal{E}$ are satisfied on every edge and only vertex configurations from $\\mathcal{V}'$ are used.\nUse the definition of $\\ell$-full set to find a coloring of any such $V_i$ and continue in a similar manner until the whole $A$ is colored.\n\\end{proof}\n\n\n\n\n\\begin{proposition}[$\\ell$-full $\\Rightarrow$ $\\mathsf{LOCAL}(O(\\log n))$]\\label{thm:ellfull_to_logn}\nLet $\\Pi = (\\Sigma, \\mathcal{V}, \\mathcal{E})$ be an LCL with an $\\ell$-full set $\\mathcal{V}' \\subseteq \\mathcal{V}$. Then $\\Pi$ can be solved in $O(\\log n)$ rounds in $\\mathsf{LOCAL}$.\n\\end{proposition}\n\\begin{proof}\nThe proof uses a variant of the rake-and-compress decomposition considered in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}.\n\n\\paragraph{The Decomposition}\nThe decomposition is parameterized an integer $\\ell' \\geq 1$, and it decomposes the vertices of $T$ into $2L - 1$ layers \\[V = \\VR{1} \\cup \\VC{1} \\cup \\VR{2} \\cup \\VC{2} \\cup \\VR{3} \\cup \\VC{3} \\cup \\cdots \\cup \\VR{L},\\] \nwith $L= O(\\log n)$. We write $\\GC{i}$ to denote the subtree induced by the vertices $\\left(\\bigcup_{j=i+1}^{L} \\VR{j}\\right) \\cup \\left( \\bigcup_{j=i}^{L-1} \\VC{j}\\right)$. Similarly, $\\GR{i}$ is the subtree induced by the vertices \n$\\left(\\bigcup_{j=i}^L \\VR{j}\\right) \\cup \\left( \\bigcup_{j=i}^{L-1} \\VC{j}\\right)$.\nThe sets $\\VR{i}$ and $\\VC{i}$ are required to satisfy the following requirements. \n\\begin{itemize}\n \\item Each $v \\in \\VR{i}$ has degree at most one in the graph $\\GR{i}$.\n \\item Each $v \\in \\VC{i}$ has degree exactly two in the graph $\\GC{i}$. Moreover, the $\\VC{i}$-vertices in $\\GC{i}$ form paths with $s$ vertices, with $\\ell' \\leq s \\leq 2 \\ell'$.\n\\end{itemize}\nFor any given constant $\\ell' \\geq 1$, it was shown in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} that such a decomposition of a tree $T$ can be computed in $O(\\log n)$ rounds. See \\cref{fig:rake_and_compress} for an example of such a decomposition with $\\ell' = 4$.\n\n\\paragraph{The Algorithm}\nGiven such a decomposition with $\\ell' = \\max\\{1, \\ell - 2\\}$, $\\Pi$ can be solved in $O(\\log n)$ rounds by labeling the vertices in this order: $\\VR{L}$, $\\VC{L-1}$, $\\VR{L-1}$, $\\ldots$, $\\VR{1}$, as follows. The algorithm only uses the vertex configurations in the $\\ell$-full set $\\mathcal{V}'$.\n\n\\paragraph{Labeling $\\VR{i}$}\nSuppose all vertices in $\\VR{L}, \\VC{L-1}, \\VR{L-1}, \\ldots, \\VC{i}$ have been labeled using $\\mathcal{V}'$. Recall that each $v \\in \\VR{i}$ has degree at most one in the graph $\\GR{i}$.\nIf $v \\in \\VR{i}$ has no neighbor in $\\VR{L} \\cup \\VC{L-1} \\cup \\VR{L-1} \\cup \\cdots \\cup \\VC{i}$, then we can label the half edges surrounding $v$ by any $c \\in \\mathcal{V}'$.\nOtherwise, $v \\in \\VR{i}$ has exactly one neighbor $u$ in $\\VR{L} \\cup \\VC{L-1} \\cup \\VR{L-1} \\cup \\cdots \\cup \\VC{i}$. Suppose the vertex configuration of $u$ is $c$, where the half-edge label on $\\{u,v\\}$ is $\\mathtt{a} \\in c$. A simple observation from the definition of $\\ell$-full sets is that for any $c \\in \\mathcal{V}'$ and any $\\mathtt{a} \\in c$, there exist $c' \\in \\mathcal{V}'$ and $\\mathtt{a}' \\in c'$ in such a way that $\\{\\mathtt{a}, \\mathtt{a}'\\} \\in \\mathcal{E}$.\nHence we can label the half edges surrounding $v$ by $c' \\in \\mathcal{V}'$ where the half-edge label on $\\{u,v\\}$ is $\\mathtt{a}' \\in c'$. \n\n\n\\paragraph{Labeling $\\VC{i}$} Suppose all vertices in $\\VR{L}, \\VC{L-1}, \\VR{L-1}, \\ldots, \\VR{i+1}$ have been labeled using $\\mathcal{V}'$. Recall that the $\\VC{i}$-vertices in $\\GC{i}$ form degree-2 paths $P=(v_1, v_2, \\ldots, v_s)$, with $\\ell' \\leq s \\leq 2 \\ell'$.\nLet $P'= (x, v_1, v_2, \\ldots, v_s, y)$ be the path resulting from appending to $P$ the neighbors of the two end-points of $P$ in $\\GC{i}$. The two vertices $x$ and $y$ are in $\\VR{L} \\cup \\VC{L-1} \\cup \\VR{L-1} \\cup \\cdots \\cup \\VR{i+1}$, so they have been assigned half-edge labels using $\\mathcal{V}'$. Since $P'$ contains at least $\\ell'+2 \\geq \\ell$ vertices, the definition of $\\ell$-full sets ensures that we can label $v_1, v_2, \\ldots, v_s$ using vertex configurations in $\\mathcal{V}'$ in such a way that the half-edge labels on $\\{x, v_1\\}, \\{v_1, v_2\\}, \\ldots, \\{v_s, y\\}$ are all in $\\mathcal{E}$.\n\\end{proof}\n\n\n\n\n\n\\begin{figure}[h!]\n \\centering\n \\includegraphics[width= .8\\textwidth]{fig\/y1.pdf}\n \\caption{The variant of the rake-and-compress decomposition used in the proof of \\cref{thm:ellfull_to_logn}. }\n \\label{fig:rake_and_compress}\n\\end{figure}\n\n\n\n\\subsection{Necessity}\n\nWe start by sketching that (2.) in \\cref{thm:MainBaireLog} is necessary for $\\mathsf{BAIRE}$ and $\\mathsf{fiid}$.\n\n\n\\begin{theorem}[Bernshteyn \\cite{Bernshteyn_work_in_progress}]\nLet $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ be an LCL and suppose that $\\Pi\\in \\mathsf{BAIRE}$ or $\\Pi\\in \\mathsf{fiid}$.\nThen $\\Pi$ admits an $\\ell$-full set $\\mathcal{V}'\\subseteq\\mathcal{V}$ for some $\\ell>0$.\n\\end{theorem}\n\\begin{proof}[Proof Sketch]\n\n\nWe start with $\\mathsf{BAIRE}$.\nSuppose that every Borel acyclic $\\Delta$-regular graph admits a Borel solution on a $\\tau$-comeager set for every compatible Polish topology $\\tau$.\nIn particular, this holds for the Borel graph induced by the standard generators of the free product of $\\Delta$-copies of $\\mathbb{Z}_2$ on the free part of the shift action on the alphabet $\\{0,1\\}$ endowed with the product topology.\nLet $F$ be such a solution.\nWrite $\\mathcal{V}'\\subseteq \\mathcal{V}$ for the configurations of half-edge labels around vertices that $F$ outputs on a non-meager set.\nLet $C$ be a comeager set on which $F$ is continuous. Then, every element of $\\mathcal{V}'$ is encoded by some finite window in the shift on $C$, that is, for each element there are a $k \\in \\mathbb{N}$ and function $s:B(1,k) \\to \\{0,1\\}$ such that $F$ is constant on the set $N_s \\cap C$ (where $N_s$ is the basic open neighbourhood determined by $s$, and $B(1,k)$ is the $k$-neighbourhood of the identity in the Cayley graph of the group).\nSince $\\mathcal{V}'$ is finite, we can take $t>0$ to be the maximum of such $k$'s.\nIt follows by standard arguments that $\\mathcal{V}'$ is $\\ell$-full for $\\ell>2t+1$.\n\nA similar argument works for the $\\mathsf{fiid}$, however, for the sake of brevity, we sketch a shorter argument that uses the fact that there must be a correlation decay for factors of iid's. Let $\\Pi\\in \\mathsf{fiid}$.\nThat is, there is an $\\operatorname{Aut}(T)$-equivariant measurable function from iid's on $T$ (without colored edges this time) into the space of $\\Pi$-colorings.\nLet $\\mathcal{V}'$ be the set of half-edges configurations around vertices that have non-zero probability to appear.\nLet $u,v\\in T$ be vertices of distance $k_0\\in \\mathbb{N}$.\nBy \\cite{BGHV} the correlation between the configurations around $u$ and $v$ tends to $0$ as $k_0\\to \\infty$.\nThis means that if the distance is big enough, then all possible pairs of $\\mathcal{V}'$ configurations need to appear.\n\\end{proof}\n\nTo finish the proof of \\cref{thm:MainBaireLog} we need to demonstrate the following theorem. Note that $\\mathsf{LOCAL}(n^{o(1)}) = \\mathsf{LOCAL}(O(\\log n))$ according to the $\\omega(\\log n)$ -- $n^{o(1)}$ complexity gap~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}.\n\n\\begin{restatable}{theorem}{logcomb}\\label{thm:logn_to_ellfull}\nLet $\\Pi = (\\Sigma, \\mathcal{V}, \\mathcal{E})$ be an LCL solvable in $\\mathsf{LOCAL}(n^{o(1)})$ rounds.\nThen there exists an $\\ell$-full set $\\mathcal{V}' \\subseteq \\mathcal{V}$ for some $\\ell \\geq 2$.\n\\end{restatable}\n\n\nThe rest of the section is devoted to the proof of \\cref{thm:logn_to_ellfull}. We start with the high-level idea of the proof.\nA natural attempt for showing $\\mathsf{LOCAL}(n^{o(1)})$ $\\Rightarrow$ $\\ell$-full is to simply take any $\\mathsf{LOCAL}(n^{o(1)})$ algorithm $\\mathcal{A}$ solving $\\Pi$, and then take $\\mathcal{V}'$ to be all vertex configurations that can possibly occur in an output of $\\mathcal{A}$. It is not hard to see that this approach does not work in general, because the algorithm might use a special strategy to label vertices with degree smaller than $\\Delta$. Specifically, there might be some vertex configuration $c$ used by $\\mathcal{A}$ so that some $\\mathtt{a} \\in c$ will only be used to label \\emph{virtual} half edges. It will be problematic to include $c$ in $\\mathcal{V}'$.\n\nTo cope with this issue, we do not deal with general bounded-degree trees. Instead, we construct recursively a sequence $(W^\\ast_1, W^\\ast_2, \\ldots, W^\\ast_L)$ of sets of rooted, layered, and \\emph{partially labeled} tree in a special manner. A tree $T$ is included in $W^\\ast_i$ if it can be constructed by gluing a multiset of rooted trees in $W^\\ast_{i-1}$ and a new root vertex $r$ in a certain fixed manner. A vertex is said to be in layer $i$ if it is introduced during the $i$-th step of the construction, i.e., it is introduced as the root $r$ during the construction of $W^\\ast_i$ from $W^\\ast_{i-1}$. All remaining vertices are said to be in layer 0.\n\n\nWe show that each $T \\in W^\\ast_L$ admits a correct labeling that extends the given partial labeling, as these partial labelings are computed by a simulation of $\\mathcal{A}$.\nMoreover, in these correct labelings, the variety of possible configurations of half-edge labels around vertices in different non-zero layers is the same for each layer.\nThis includes vertices of non-zero layer whose half-edges are labeled by the given partial labeling.\nWe simply pick $\\mathcal{V}'$ to be the set of all configurations of half-edge labels around vertices that can appear in a non-zero layer in a correct labeling of a tree $T \\in W^\\ast_L$.\nOur construction ensures that each $c \\in \\mathcal{V}'$ appears as the labeling of some degree-$\\Delta$ vertex in some tree that we consider.\n\n\nThe proof that $\\mathcal{V}'$ is an $\\ell$-full set is based on finding paths using vertices of non-zero layers connecting two vertices with any two vertex configurations in $\\mathcal{V}'$ in different lengths. These paths exist because the way rooted trees in $W^\\ast_{i-1}$ are glued together in the construction of $W^\\ast_{i}$ is sufficiently flexible.\nThe reason that we need $\\mathcal{A}$ to have complexity $\\mathsf{LOCAL}(n^{o(1)})$ is that the construction of the trees can be parameterized by a number $w$ so that all the trees have size polynomial in $w$ and the vertices needed to be assigned labeling are at least distance $w$ apart from each other. Since the number of rounds of $\\mathcal{A}$ executed on trees of size $w^{O(1)}$ is much less than $w$, each labeling assignment can be calculated locally and independently. \nThe construction of the trees as well as the analysis are based on a machinery developed in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}.\n Specifically, we will consider the equivalence relation $\\overset{\\star}{\\sim}$ defined in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} and prove some of its properties, including a pumping lemma for bipolar trees. The exact definition of $\\overset{\\star}{\\sim}$ in this paper is different from the one in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} because the mathematical formalism describing LCL problems in this paper is different from the one in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}.\nAfter that, we will consider a procedure for gluing trees parameterized by a labeling function $f$ similar to the one used in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}. \nWe will apply this procedure iteratively to generate a set of trees. We will show that the desired $\\ell$-full set $\\mathcal{V}' \\subseteq \\mathcal{V}$ can be constructed by considering the set of all possible correct labeling of these trees. \n \n\n\\paragraph{The Equivalence Relation $\\overset{\\star}{\\sim}$} \nWe consider trees with a list of designated vertices $v_1, v_2, \\ldots, v_k$ called \\emph{poles}. A \\emph{rooted tree} is a tree with one pole $r$, and a \\emph{bipolar tree} is a tree with two poles $s$ and $t$. \nFor a tree $T$ with its poles $S=(v_1, v_2, \\ldots, v_k)$ with $\\deg(v_i) = d_i < \\Delta$, we denote by $h_{(T,S)}$ the function that maps each choice of the \\emph{virtual} half-edge labeling surrounding the poles of $T$ to $\\textsf{YES}$ or $\\textsf{NO}$, indicating whether such a partial labeling can be completed into a correct complete labeling of $T$. \nMore specifically, consider \n\\[X=(\\mathcal{I}_1, \\mathcal{I}_2, \\ldots, \\mathcal{I}_k),\\]\nwhere $\\mathcal{I}_i$ is a size-$(\\Delta - d_i)$ multiset of labels in $\\Sigma$, for each $1 \\leq i \\leq k$. Then $h_{(T,S)}(X) = \\textsf{YES}$ if there is a correct labeling of $T$ such that the $\\Delta - d_i$ virtual half-edge labels surrounding $v_i$ are labeled by $\\mathcal{I}_i$, for each $1 \\leq i \\leq k$. This definition can be generalized to the case $T$ is already partially labeled in the sense that some of the half-edge labels have been fixed. In this case, $h_{(T,S)}(X) = \\textsf{NO}$ whenever $X$ is incompatible with the given partial labeling.\n\nLet $T$ be a tree with poles $S = (v_1, v_2, \\ldots, v_k)$ and let $T'$ be another tree with poles $S' = (v_1', v_2', \\ldots, v_k')$ such that $\\deg(v_i) = \\deg(v_i') = d_i$ for each $1 \\leq i \\leq k$.\nThen we write $T_1 \\overset{\\star}{\\sim} T_2$ if $h_{(T,S)} = h_{(T',S')}$.\n\nGiven an LCL problem $\\Pi$, it is clear that the number of equivalence classes of rooted trees and bipolar trees w.r.t.~$\\overset{\\star}{\\sim}$ is finite.\nFor a rooted tree $T$, denote by $\\textsf{Class}_1(T)$ the equivalence class of $T$. For a bipolar tree $H$, denote by $\\texttt{type}(H)$ the equivalence class of $H$.\n\n\n\n\n\\paragraph{Subtree Replacement}\nThe following lemma provides a sufficient condition that the equivalence class of a tree $T$ is invariant of the equivalence class of its subtree $T'$.\nWe note that a real half edge in $T$ might become virtual in its subtree $T'$.\nConsider a vertex $v$ in $T'$ and its neighbor $u$ that is in $T$ but not in $T'$. Then the half edge $(v, \\{u,v\\})$ is real in $T$ and virtual in $T'$.\n\n\n\\begin{lemma}[Replacing subtrees]\\label{lem:replace}\nLet $T$ be a tree with poles $S$.\nLet $T'$ be a connected subtree of $T$ induced by $U \\subseteq V$, where $V$ is the set of vertices in $T$.\nWe identify a list of designated vertices $S' = (v_1', v_2', \\ldots, v_k')$ in $U$ satisfying the following two conditions to be the poles of $T'$.\n\\begin{itemize}\n \\item $S \\cap U \\subseteq S'$.\n \\item Each edge $e=\\{u,v\\}$ connecting $u \\in U$ and $ v \\in V \\setminus U$ must satisfy $u \\in S'$.\n\\end{itemize}\nLet $T''$ be another tree with poles $S'' = (v_1'', v_2'', \\ldots, v_k'')$ that is in the same equivalence class as $T'$.\nLet $T^\\ast$ be the result of replacing $T'$ by $T''$ in $T$ by identifying $v_i' = v_i''$ for each $1 \\leq i \\leq k$.\nThen $T^\\ast$ is in the same equivalence class as $T$. \n\\end{lemma}\n\\begin{proof}\nTo prove the lemma, by symmetry, it suffices to show that starting from any correct labeling $\\mathcal{L}$ of $T$, it is possible to find a correct labeling $\\mathcal{L}^\\ast$ of $T^\\ast$ in such a way that the multiset of the\nvirtual half-edge labels surrounding each pole in $S$ remain the same.\n\nSuch a correct labeling $\\mathcal{L}^\\ast$ of $T^\\ast$ is constructed as follows. If $v \\in V \\setminus U$, then we simply adopt the given labeling $\\mathcal{L}$ of $v$ in $T$.\nNext, consider the vertices in $U$.\nSet $X=(\\mathcal{I}_1, \\mathcal{I}_2, \\ldots, \\mathcal{I}_k)$ to be the one compatible with the labeling $\\mathcal{L}$ of $T$ restricted to the subtree $T'$ in the sense that $\\mathcal{I}_i$ is the multiset of the virtual half-edge labels surrounding the pole $v_i'$ of $T'$, for each $1 \\leq i \\leq k$. We must have $h_{T',S'}(X) = \\textsf{YES}$.\nSince $T'$ and $T''$ are in the same equivalence class, we have $h_{T'',S''}(X) = \\textsf{YES}$, and so we can find a correct labeling $\\mathcal{L}''$ of $T''$ that is also compatible with $X$. Combining this labeling $\\mathcal{L}''$ of the vertices $U$ in $T''$ with the labeling $\\mathcal{L}$ of the vertices $V \\setminus U$ in $T$, we obtain a desired labeling $\\mathcal{L}^\\ast$ of $T^\\ast$.\n\nWe verify that $\\mathcal{L}^\\ast$ gives a correct labeling of $T^\\ast$. Clearly, the size-$\\Delta$ multiset that labels each vertex $v \\in V$ is in $\\mathcal{V}$, by the correctness of $\\mathcal{L}$ and $\\mathcal{L}''$. Consider any edge $e=\\{u,v\\}$ in $T^\\ast$. Similarly, if $\\{u,v\\} \\subseteq V \\setminus U$ or $\\{u,v\\} \\cap (V \\setminus U) = \\emptyset$, then the size-$2$ multiset that labels $e$ is in $\\mathcal{E}$, by the correctness of $\\mathcal{L}$ and $\\mathcal{L}''$. For the case that $e=\\{u,v\\}$ connects a vertex $u \\notin V$ and a vertex $v \\in V \\setminus U$, we must have $u \\in S'$ by the lemma statement. Therefore, the label of the half edge $(u,e)$ is the same in both $\\mathcal{L}$ and $\\mathcal{L}^\\ast$ by our choice of $\\mathcal{L}''$. Thus, the size-$2$ multiset that labels $e$ is in $\\mathcal{E}$, by the correctness of $\\mathcal{L}$.\n\nWe verify that the \nvirtual half-edge labels surrounding each pole in $S$ are the same in both $\\mathcal{L}$ and $\\mathcal{L}^\\ast$. Consider a pole $v \\in S$. If $v \\in V \\setminus U$, then the labeling of $v$ is clearly the same in both $\\mathcal{L}$ and $\\mathcal{L}^\\ast$.\nIf $v \\notin V\\setminus U$, then the condition $S \\cap U \\subseteq S'$ in the statement implies that $v \\in S'$. In this case, the way we pick $\\mathcal{L}''$ ensures that the \nvirtual half-edge labels surrounding $v \\in S'$ are the same in both $\\mathcal{L}$ and $\\mathcal{L}^\\ast$.\n\\end{proof}\n\nIn view of the proof of \\cref{lem:replace}, as long as the conditions in \\cref{lem:replace} are met, we are able to abstract out a subtree by its equivalence class when reasoning about correct labelings of a tree. This observation will be applied repeatedly in the subsequent discussion.\n\n\n\\paragraph{A Pumping Lemma} We will prove a pumping lemma of bipolar trees using \\cref{lem:replace}. \nSuppose $T_i$ is a tree with a root $r_i$ for each $1 \\leq i \\leq k$, then $H=(T_1, T_2, \\ldots, T_k)$ denotes the bipolar tree resulting from concatenating the roots $r_1, r_2, \\ldots, r_k$ into a path $(r_1, r_2, \\ldots, r_k)$ and setting the two poles of $H$ by $s = r_1$ and $t = r_k$.\nA simple consequence of \\cref{lem:replace} is that $\\texttt{type}(H)$ is determined by $\\textsf{Class}_1(T_1), \\textsf{Class}_1(T_2), \\ldots, \\textsf{Class}_1(T_k)$.\nWe have the following pumping lemma. \n\n\\begin{lemma}[Pumping lemma]\\label{lem:pump}\nThere exists a finite number $ {\\ell_{\\operatorname{pump}}} > 0$ such that as long as $k \\geq {\\ell_{\\operatorname{pump}}} $, any bipolar tree $H=(T_1, T_2, \\ldots, T_k)$ can be decomposed into $H = X \\circ Y \\circ Z$ with $0 < |Y| < k$ so that $X \\circ Y^i \\circ Z$ is in the same equivalence class as $H$ for each $i \\geq 0$. \n\\end{lemma}\n\\begin{proof}\nSet $ {\\ell_{\\operatorname{pump}}} $ to be the number of equivalence classes for bipolar trees plus one. By the pigeon hole principle, there exist $1 \\leq a < b \\leq k$ such that $(T_1, T_2, \\ldots, T_a)$ and $(T_1, T_2, \\ldots, T_b)$ are in the same equivalence class. Set $X = (T_1, T_2, \\ldots, T_a)$, $Y = (T_{a+1}, T_{a+2}, \\ldots, T_b)$, and $Z = (T_{b+1}, T_{b+2}, \\ldots, T_k)$. As we already know that $\\texttt{type}(X) = \\texttt{type}(X\\circ Y)$, \\cref{lem:replace} implies that \n$\\texttt{type}(X\\circ Y) = \\texttt{type}(X^2\\circ Y)$ by replacing $X$ by $X\\circ Y$ in the bipolar tree $X\\circ Y$.\nSimilarly, $\\texttt{type}(X \\circ Y^i)$ is the same for for each $i \\geq 0$. Applying \\cref{lem:replace} again to replace $X \\circ Y$ by $X \\circ Y^i$ in $H = X \\circ Y \\circ Z$, we conclude that $X \\circ Y^i \\circ Z$ is in the same equivalence class as $H$ for each $i \\geq 0$. \n\\end{proof}\n\n\n\n\\paragraph{A Procedure for Gluing Trees} Suppose that we have a set of rooted trees $W$. We devise a procedure that generates a new set of rooted trees by gluing the rooted trees in $W$ together. \nThis procedure is parameterized by a labeling function $f$. Consider a bipolar tree \\[H = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)\\] \nwhere $T^m$ is formed by attaching the roots $r_1, r_2, \\ldots, r_{\\Delta - 2}$ of the rooted trees $T_1^m, T_2^m, \\ldots, T_{\\Delta -2}^m$ to the root $r^m$ of $T^m$.\n\n\nThe labeling function $f$ assigns the half-edge labels surrounding $r^m$ based on \n\\[\\texttt{type}(H^l), \\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m), \\texttt{type}(H^r)\\]\nwhere $H^l = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l)$ and $H^r = (T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)$.\n\nWe write $T^{m}_\\ast$ to denote the result of applying $f$ to label the root $r^m$ of $T^{m}$ in the bipolar tree $H$, and we write \\[H_\\ast = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T_\\ast^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)\\] to denote the result of applying $f$ to $H$. We make the following observation.\n\n\\begin{lemma}[Property of $f$]\\label{lem:labeling_function}\nThe two equivalence classes\n$\\textsf{Class}_1(T_\\ast^m)$ and $\\texttt{type}(H_\\ast)$ are determined by \\[\\texttt{type}(H^l), \\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m), \\texttt{type}(H^r),\\] and the labeling function $f$.\n\\end{lemma}\n\\begin{proof}\nBy \\cref{lem:replace}, once the half-edge labelings of the root $r^m$ of $T^{m}$ is fixed, $\\textsf{Class}_1(T_\\ast^m)$ is determined by $\\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m)$. \nTherefore, indeed $\\textsf{Class}_1(T_\\ast^m)$ is determined by\n\\[\\texttt{type}(H^l), \\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m), \\texttt{type}(H^r),\\] and the labeling function $f$.\nSimilarly, applying \\cref{lem:replace} to the decomposition of $H_\\ast$ into $H^l$, $T_\\ast^m$, and $H^r$, we infer that $\\texttt{type}(H_\\ast)$ depends only on $\\texttt{type}(H^l)$, $\\textsf{Class}_1(T_\\ast^m)$, and $\\texttt{type}(H^r)$.\n\\end{proof}\n\nThe three sets of trees $X_f(W)$, $Y_f(W)$, and $Z_f(W)$ are constructed as follows.\n\n\\begin{itemize}\n \\item $X_f(W)$ is the set of all rooted trees resulting from appending $\\Delta-2$ arbitrary rooted trees $T_1, T_2, \\ldots, T_{\\Delta-2}$ in $W$ to a new root vertex $r$.\n \\item $Y_f(W)$ is the set of bipolar trees constructed as follows. For each choice of $2 {\\ell_{\\operatorname{pump}}} +1$ rooted trees $T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r$ from $X_f(W)$, concatenate them into a bipolar tree \n \\[H = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r),\\] let $H_\\ast$ be the result of applying the labeling function $f$ to $H$, and then add $H_\\ast$ to $Y_f(W)$.\n \\item $Z_f(W)$ is the set of rooted trees constructed as follows. For each \\[H_\\ast = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T_\\ast^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r) \\in Y_f(W),\\] add $(T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T_\\ast^m)$ to $Z_f(W)$, where we set the root of $T_1^l$ as the root, and add $(T_\\ast^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)$ to $Z_f(W)$, where we set the root of ${T_{ {\\ell_{\\operatorname{pump}}} }^r}$ as the root. \n\\end{itemize}\n\nWe write $\\textsf{Class}_1(S) = \\bigcup_{T \\in S} \\{\\textsf{Class}_1(T)\\}$ and $\\texttt{type}(S) = \\bigcup_{H \\in S} \\{\\texttt{type}(H)\\}$, and \nwe make the following observation.\n\n\\begin{lemma}[Property of $X_f(W)$, $Y_f(W)$, and $Z_f(W)$]\\label{lem:sets_of_trees}\nThe three sets of equivalence classes $\\textsf{Class}_1(X_f(W))$, $\\texttt{type}(Y_f(W))$, and $\\textsf{Class}_1(Z_f(W))$ depend only on $\\textsf{Class}_1(W)$ and the labeling function $f$.\n\\end{lemma}\n\\begin{proof}\nThis is a simple consequence of \\cref{lem:replace,lem:labeling_function}.\n\\end{proof}\n\n\\paragraph{A Fixed Point $W^\\ast$}\nGiven a fixed labeling function $f$,\nwe want to find a set of rooted trees $W^\\ast$ that is a \\emph{fixed point} for the procedure $Z_f$ in the sense that \\[\\textsf{Class}_1(Z_f(W^\\ast)) = \\textsf{Class}_1(W^\\ast).\\] \nTo find such a set $W^\\ast$, we construct a two-dimensional array of rooted trees $\\{W_{i,j}\\}$, as follows.\n\n\\begin{itemize}\n \\item For the base case, $W_{1,1}$ consists of only the one-vertex rooted tree.\n \\item Given that $W_{i,j}$ as been constructed, we define $W_{i,j+1} = Z_f(W_{i,j})$.\n \\item Given that $W_{i,j}$ for all positive integers $j$ have been constructed, $W_{i+1,1}$ is defined as follows. Pick $b_i$ as the smallest index such that $\\textsf{Class}_1(W_{i, b_i}) = \\textsf{Class}_1(W_{i, a_i})$ for some $1 \\leq a_i < b_i$. By the pigeon hole principle, the index $b_i$ exists, and it is upper bounded by $2^C$, where $C$ is the number of equivalence classes for rooted trees. We set \n \\[W_{i+1,1} = W_{i, a_i} \\cup W_{i, a_i + 1} \\cup \\cdots \\cup W_{i, b_i - 1}.\\]\n\\end{itemize}\n\nWe show that the sequence $\\textsf{Class}_1(W_{i,a_i}), \\textsf{Class}_1(W_{i,a_i+1}), \\textsf{Class}_1(W_{i,a_i+2}), \\ldots$ is periodic with a period $c_i = b_i - a_i$.\n\n\\begin{lemma}\\label{lem:W-aux-1}\nFor any $i \\geq 1$ and for any $j \\geq a_i$, we have $\\textsf{Class}_1(W_{i,j}) = \\textsf{Class}_1(W_{i,j+c_i})$, where $c_i = b_i - a_i$.\n\\end{lemma}\n\\begin{proof}\nBy \\cref{lem:sets_of_trees}, $\\textsf{Class}_1(W_{i,j})$ depends only on $\\textsf{Class}_1(W_{i,j-1})$. Hence the lemma follows from the fact that $\\textsf{Class}_1(W_{i, b_i}) = \\textsf{Class}_1(W_{i, a_i})$.\n\\end{proof}\n\nNext, we show that $\\textsf{Class}_1(W_{2,1}) \\subseteq \\textsf{Class}_1(W_{3, 1}) \\subseteq \\textsf{Class}_1(W_{4, 1}) \\subseteq \\cdots$.\n\n\\begin{lemma}\\label{lem:W-aux-2}\nFor any $i \\geq 2$, we have $\\textsf{Class}_1(W_{i,1}) \\subseteq \\textsf{Class}_1(W_{i,j})$ for each $j > 1$, and so $\\textsf{Class}_1(W_{i,1}) \\subseteq \\textsf{Class}_1(W_{i+1, 1})$.\n\\end{lemma}\n\\begin{proof}\nSince $W_{i,1} = \\bigcup_{a_{i-1} \\leq l \\leq b_{i-1} - 1} W_{i-1, l}$, we have \\[\\bigcup_{a_{i-1}+j-1 \\leq l \\leq b_{i-1} +j} W_{i-1, l} \\subseteq W_{i,j}\\] according to the procedure of constructing $W_{i,j}$. By \\cref{lem:W-aux-1}, \\[ \\textsf{Class}_1\\left(\\bigcup_{a_{i-1} \\leq l \\leq b_{i-1} - 1} W_{i-1, l}\\right) = \\textsf{Class}_1\\left(\\bigcup_{a_{i-1}+j-1 \\leq l \\leq b_{i-1} +j} W_{i-1, l}\\right)\\] for all $j \\geq 1$, and so we have $\\textsf{Class}_1(W_{i,1}) \\subseteq \\textsf{Class}_1(W_{i,j})$ for each $j > 1$. The claim $\\textsf{Class}_1(W_{i,1}) \\subseteq \\textsf{Class}_1(W_{i+1, 1})$ follows from the fact that $W_{i+1,1} = \\bigcup_{a_{i} \\leq l \\leq b_{i} - 1} W_{i, l}$.\n\\end{proof}\n\nSet $i^\\ast$ to be the smallest index $i \\geq 2$ such that $\\textsf{Class}_1(W_{i,1}) = \\textsf{Class}_1(W_{i+1, 1})$. By the pigeon hole principle and \\cref{lem:W-aux-2}, the index $i^\\ast$ exists, and it is upper bounded by $C$, the number of equivalence classes for rooted trees. We set \\[W^\\ast = W_{i^\\ast, 1}.\\] \n\nThe following lemma shows that $\\textsf{Class}_1(Z_f(W^\\ast)) = \\textsf{Class}_1(W^\\ast)$, as needed.\n\n\\begin{lemma}\\label{lem:W-aux-3}\nFor any $i \\geq 2$, if $\\textsf{Class}_1(W_{i,1}) = \\textsf{Class}_1(W_{i+1, 1})$, then $\\textsf{Class}_1(W_{i,j})$ is the same for all $j \\geq 1$. \n\\end{lemma}\n\\begin{proof}\nBy \\cref{lem:W-aux-1} and the way we construct $W_{i+1, 1}$, we have \\[\\textsf{Class}_1(W_{i, j}) \\subseteq \\textsf{Class}_1(W_{i+1, 1}) \\ \\ \\text{for each} \\ \\ j \\geq a_i.\\]\nBy \\cref{lem:W-aux-2}, we have \\[\\textsf{Class}_1(W_{i, 1}) \\subseteq \\textsf{Class}_1(W_{i, j}) \\ \\ \\text{for each} \\ \\ j > 1.\\]\nTherefore, $\\textsf{Class}_1(W_{i,1}) = \\textsf{Class}_1(W_{i+1, 1})$ implies that $\\textsf{Class}_1(W_{i, j}) = \\textsf{Class}_1(W_{i,1})$ for each $j \\geq a_i$. Hence we must have $\\textsf{Class}_1(Z_f(W_{i,1})) = \\textsf{Class}_1(W_{i,1})$, and so $\\textsf{Class}_1(W_{i,j})$ is the same for all $j \\geq 1$. \n\\end{proof}\n\nWe remark that simply selecting $W^\\ast$ to be any set such that $\\textsf{Class}_1(Z_f(W^\\ast)) = \\textsf{Class}_1(W^\\ast)$ is not enough for our purpose. As we will later see, it is crucial that the set $W^\\ast$ is constructed by iteratively applying the function $Z_f$ and taking the union of previously constructed sets.\n\n\n\n\n\n\n\n\\paragraph{A Sequence of Sets of Trees}\nWe define $W_1^\\ast = W^\\ast$ and $W_i^\\ast = Z_f(W_{i-1}^\\ast)$ for each $1 < i \\leq L$, where $L$ is some sufficiently large number to be determined.\n\nThe way we choose $W^\\ast$ guarantees that $\\textsf{Class}_1(W_i^\\ast) = \\textsf{Class}_1(W^\\ast)$ for all $1 \\leq i \\leq L$. \nFor convenience, we write $X_i^\\ast = X_f(W_i^\\ast)$ and $Y_i^\\ast = Y_f(W_i^\\ast)$.\nSimilarly, $\\textsf{Class}_1(X_i^\\ast)$ is the same for all $1 \\leq i \\leq L$ and $\\texttt{type}(Y_i^\\ast)$ is the same for all $1 \\leq i \\leq L$.\n\n\nOur analysis will rely on the assumption that all rooted trees in $W^\\ast$ admit correct labelings. Whether this is true depends only on $\\textsf{Class}_1(W^\\ast)$, which depends only on the labeling function $f$. We say that $f$ is \\emph{feasible} if it leads to a set $W^\\ast$ where all the rooted trees therein admit correct labelings. The proof that a feasible labeling function $f$ exists is deferred.\n\nThe assumption that $f$ is feasible implies that all trees in $W_i^\\ast$, $X_i^\\ast$, and $Y_i^\\ast$, for all $1 \\leq i \\leq L$, admit correct labelings. All rooted trees in $X_i^\\ast$ admit correct labelings because they are subtrees of the rooted trees in $W_{i+1}^\\ast$. A correct labeling of any bipolar tree $H \\in Y_i^\\ast$ can be obtained by combining any correct labelings of the two rooted trees in $W_{i+1}^\\ast$ resulting from $H$.\n\n\n\n\n\\paragraph{Layers of Vertices}\nWe assign a layer number $\\lambda(v)$ to each vertex $v$ in a tree based on the step that $v$ is introduced in the construction \\[W_1^\\ast \\rightarrow W_2^\\ast \\rightarrow \\cdots \\rightarrow W_L^\\ast.\\]\n\nIf a vertex $v$ is introduced as the root vertex of a tree in $X_i^\\ast = X_f(W_i^\\ast)$, then we say that the layer number of $v$ is $\\lambda(v) = i \\in \\{1,2, \\ldots, L\\}$. A vertex $v$ has $\\lambda(v) = 0$ if it belongs to a tree in $W_1^\\ast$.\n\n\nFor any vertex $v$ with $\\lambda(v) = i$ in a tree $T \\in X_j^\\ast$ with $i \\leq j$, we write $T_v$ to denote the subtree of $T$ such that $T_v \\in X_i^\\ast$ where $v$ is the root of $T_v$. \n\n\n \nWe construct a sequence of sets $R_1, R_2, \\ldots, R_L$ as follows.\nWe go over all rooted trees $T \\in X_L^\\ast$, all possible correct labeling $\\mathcal{L}$ of $T$, and all vertices $v$ in $T$ with $\\lambda(v) = i \\in \\{1,2, \\ldots, L\\}$. Suppose that the $\\Delta-2$ rooted trees in the construction of $T_v \\in X_i^\\ast = X_f(W_i^\\ast)$ are $T_1$, $T_2$, $\\ldots$, $T_{\\Delta-2}$, and let $r_i$ be the root of $T_i$. Consider the following parameters.\n\\begin{itemize}\n \\item $c_i = \\textsf{Class}_1(T_i)$.\n \\item $\\mathtt{a}_i$ is the real half-edge label of $v$ in $T_v$ for the edge $\\{v, r_i\\}$, under the correct labeling $\\mathcal{L}$ of $T$ restricted to $T_v$.\n \\item $\\mathcal{I}$ is the size-$2$ multiset of virtual half-edge labels of $v$ in $T_v$, under the correct labeling $\\mathcal{L}$ of $T$ restricted to $T_v$.\n\\end{itemize}\nThen we add \n$(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I})$ to $R_i$.\n\n\n\n\n\\begin{lemma}\\label{lem:Rsets1}\nFor each $1 \\leq i < L$, $R_i$ is determined by $R_{i+1}$.\n\\end{lemma}\n\\begin{proof}\nWe consider the following alternative way of constructing $R_i$ from $R_{i+1}$.\nEach rooted tree $T'$ in $W_{i+1}^\\ast = Z_f(W_i^\\ast)$ can be described as follows.\n\\begin{itemize}\n \\item Start with a path $(r_1, r_2, \\ldots, r_{ {\\ell_{\\operatorname{pump}}} + 1})$, where $r_1$ is the root of $T'$.\n \\item For each $1 \\leq j \\leq {\\ell_{\\operatorname{pump}}} +1$, append $\\Delta-2$ rooted trees $T_{j,1}, T_{j,2}, \\ldots, T_{j, \\Delta-2} \\in W_i^\\ast$ to $r_j$.\n \\item Assign the labels to the half edges surrounding $r_{ {\\ell_{\\operatorname{pump}}} + 1}$ according to the labeling function $f$.\n\\end{itemize}\n\n Now, consider the function $\\phi$ that maps each equivalence class $c$ for rooted trees to a subset of $\\Sigma$ defined as follows:\n$\\mathtt{a} \\in \\phi(c)$ if there exist \\[(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I}) \\in R_{i+1}\\] and $1 \\leq j \\leq \\Delta -2$ such that $c = c_j$ and $\\mathtt{a} = \\mathtt{a}_j$.\n\nWe go over all possible $T' \\in W_{i+1}^\\ast = Z_f(W_i^\\ast)$.\nNote that the root $r$ of $T'$ has exactly one virtual half edge. For each $\\mathtt{b} \\in \\Sigma$ such that $\\{\\mathtt{a}, \\mathtt{b}\\} \\in \\mathcal{E}$ for some $\\mathtt{a} \\in \\phi(\\textsf{Class}_1(T'))$, we go over all possible correct labelings $\\mathcal{L}$ of $T'$ where the virtual half edge of $r$ is labeled $\\mathtt{b}$. For each $1 \\leq j \\leq {\\ell_{\\operatorname{pump}}} +1$, consider the following parameters.\n\\begin{itemize}\n \\item $c_l = \\textsf{Class}_1(T_{j,l})$.\n \\item $\\mathtt{a}_l$ is the half-edge label of $r_j$ for the edge $\\{r_j, r_{j,l}\\}$.\n \\item $\\mathcal{I}$ is the size-$2$ multiset of the remaining two half-edge labels of $v$.\n\\end{itemize}\nThen we add \n$(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I})$ to $R_i$. \n\nThis construction of $R_i$ is equivalent to the original construction of $R_i$ because a correct labeling of $T' \\in W_i^\\ast$ can be extended to a correct labeling of a tree $T \\in X_L^\\ast$ that contains $T'$ as a subtree if and only if the virtual half-edge label of the root $r$ of $T'$ is $\\mathtt{b} \\in \\Sigma$ such that $\\{\\mathtt{a}, \\mathtt{b}\\} \\in \\mathcal{E}$ for some $\\mathtt{a} \\in \\phi(\\textsf{Class}_1(T'))$.\n\n\nIt is clear that this construction of $R_i$ only depends on $\\textsf{Class}_1(W_i^\\ast)$, the labeling function $f$, and the function $\\phi$, which depends only on $R_{i+1}$. Since the labeling function $f$ is fixed and $\\textsf{Class}_1(W_i^\\ast)$ is the same for all $i$, we conclude that $R_{i}$ depends only on $R_{i+1}$.\n\\end{proof}\n\n\\begin{lemma}\\label{lem:Rsets2}\nWe have $R_1 \\subseteq R_2 \\subseteq \\cdots \\subseteq R_L$.\n\\end{lemma}\n\\begin{proof}\nFor the base case, we show that $R_{L-1} \\subseteq R_L$. In fact, our proof will show that $R_{i} \\subseteq R_L$ for each $1 \\leq i < L$. Consider any \n\\[(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I}) \\in R_i.\\] \nThen there is a rooted tree $T \\in X_i^\\ast$ that is formed by attaching $\\Delta-2$ rooted trees of equivalence classes $c_1, c_2, \\ldots, c_{\\Delta-2}$ to the root vertex so that if we label the half edges surrounding the root vertex according to $\\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}$, and $\\mathcal{I}$, then this partial labeling can be completed into a correct labeling of $T$.\n\nBecause $\\textsf{Class}_1(W_i^\\ast) = \\textsf{Class}_1(W_L^\\ast)$, there is also a rooted tree $T' \\in X_L^\\ast$ that is formed by attaching $\\Delta-2$ rooted trees of equivalence classes $c_1, c_2, \\ldots, c_{\\Delta-2}$ to the root vertex. Therefore, if we label the root vertex of $T'$ in the same way as we do for $T$, then this partial labeling can also be completed into a correct labeling of $T'$. Hence we must have \n\\[(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I}) \\in R_L.\\] \n\n\nNow, suppose that we already have $R_i \\subseteq R_{i+1}$ for some $1 < i < L$. We will show that $R_{i-1} \\subseteq R_{i}$. \nDenote by $\\phi_{i}$ and $\\phi_{i+1}$ the function $\\phi$ in \\cref{lem:Rsets1} constructed from $R_{i}$ and $R_{i+1}$. \nWe have $\\phi_i(c) \\subseteq \\phi_{i+1}(c)$ for each equivalence class $c$, because $R_i \\subseteq R_{i+1}$. Therefore, in view of the alternative construction described in the proof of \\cref{lem:Rsets1}, we have $R_{i-1} \\subseteq R_{i}$.\n\\end{proof}\n\n\n\n\\paragraph{The Set of Vertex Configurations $\\mathcal{V}'$}\nBy \\cref{lem:Rsets1,lem:Rsets2}, if we pick $L$ to be sufficient large, we can have $R_1 = R_2 = R_3$.\nMore specifically, if we pick \n\\[L \\geq C^{\\Delta-2} \\cdot \\binom{|\\Sigma|+1}{|\\Sigma|-1} + 3,\\]\nthen there exists an index $3 \\leq i \\leq L$ such that $R_i = R_{i-1}$, implying that $R_1 = R_2 = R_3$. Here $C$ is the number of equivalence classes for rooted trees and $\\binom{|\\Sigma|+1}{|\\Sigma|-1}$ is the number of size-$2$ multisets of elements from $\\Sigma$.\n\n\n\nThe set $\\mathcal{V}'$ is defined by including all size-$\\Delta$ multisets $\\{\\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}\\} \\cup \\mathcal{I}$ such that \\[(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I}) \\in R_{1}\\] for some $c_1, c_2, \\ldots, c_{\\Delta-2}$.\n\n\n\n\\paragraph{The Set $\\mathcal{V}'$ is $\\ell$-full}\nTo show that the set $\\mathcal{V}'$ is $\\ell$-full, we consider the subset $\\mathcal{V}^\\ast \\subseteq \\mathcal{V}'$ defined by the set of vertex configurations used by the labeling function $f$ in the construction of $Y_f(W_i^\\ast)$. The definition of $\\mathcal{V}^\\ast$ is invariant of $i$ as $\\textsf{Class}_1(W_i^\\ast)$ is the same for all $i$. Clearly, for each $x \\in \\mathcal{V}^\\ast$, there is a rooted tree $T \\in W_2^\\ast$ where the root of a subtree $T' \\in X_1^\\ast$ of $T$ has its size-$\\Delta$ multiset of half-edge labels fixed to be $x$ by $f$, and so $\\mathcal{V}^\\ast \\subseteq \\mathcal{V}'$.\n\nFor notational simplicity, we write $x \\overset{k}{\\leftrightarrow} x'$ if there is a correct labeling of a $k$-vertex path $(v_1, v_2, \\ldots, v_k)$ using only vertex configurations in $\\mathcal{V}'$ so that the vertex configuration of $v_1$ is $x$ and the vertex configuration of $v_k$ is $x'$. If it is further required that the half-edge label of $v_1$ for the edge $\\{v_1, v_2\\}$ is $\\mathtt{a} \\in x$, then we write $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} x'$. The notation $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{a}')$ is defined similarly.\n\n\n\\begin{lemma}\\label{lem:find_path_1}\nFor any $x \\in \\mathcal{V}' \\setminus \\mathcal{V}^\\ast$ and $\\mathtt{a} \\in x$, there exist a vertex configuration $x' \\in \\mathcal{V}^\\ast$ and a number $2 \\leq k \\leq 2 {\\ell_{\\operatorname{pump}}} + 1$ such that $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} x'$.\n\\end{lemma}\n\\begin{proof}\nLet $T \\in W_L^\\ast$ be chosen so that there is a correct labeling $\\mathcal{L}$ where $x$ is a vertex configuration of some vertex $v$ with $\\lambda(v) = 2$.\n\nTo prove the lemma, it suffices to show that for each of the $\\Delta$ neighbors $u$ of $v$, it is possible to find a path $P = (v,u, \\ldots, w)$ meeting the following conditions.\n\\begin{itemize}\n \\item $w$ is a vertex whose half-edge labels have been fixed by $f$. This ensures that the vertex configuration of $w$ is in $\\mathcal{V}^\\ast$.\n \\item All vertices in $P$ are within layers 1,2, and 3. This ensures that the vertex configuration of all vertices in $P$ are in $\\mathcal{V}'$.\n \\item The number of vertices $k$ in $P$ satisfies $2 \\leq k \\leq 2 {\\ell_{\\operatorname{pump}}} + 1$. \n\\end{itemize}\n\nWe divide the proof into three cases. Refer to \\cref{fig:ell_full_construction1} for an illustration, where squares are vertices whose half-edge labels have been fixed by $f$.\n\n\\paragraph{Case 1} Consider the subtree $T_v \\in X_2^\\ast$ of $T$ whose root is $v$. In view of the construction of the set $X_f(W_2^\\ast)$, $v$ has $\\Delta-2$ children $u_1, u_2, \\ldots, u_{\\Delta-2}$ in $T_v$, where the subtree $T_i$ rooted at $u_i$ is a rooted tree in $W_2^\\ast$. \n\nFor each $1 \\leq i \\leq \\Delta-2$, according to the structure of the trees in the set $W_2^\\ast = Z_f(W_1^\\ast)$, there is a path $(u_i = w_1, w_2, \\ldots, w_{ {\\ell_{\\operatorname{pump}}} +1})$ in $T_i$ containing only layer-1 vertices, where the half-edge labels of $w_{ {\\ell_{\\operatorname{pump}}} +1}$ have been fixed by $f$. Hence $P = (v, u_i = w_1, w_2, \\ldots, w_{ {\\ell_{\\operatorname{pump}}} +1})$ is a desired path with $k = {\\ell_{\\operatorname{pump}}} + 2$ vertices.\n \n \n \n \n \n\\paragraph{Case 2} Consider the subtree $T' \\in W_3^\\ast = Z_f(W_2^\\ast)$ that contains $v$ in $T$.\nSimilarly, according to the structure of the trees in the set $W_3^\\ast = Z_f(W_2^\\ast)$, there is a path $(r = w_1', w_2', \\ldots, w_{ {\\ell_{\\operatorname{pump}}} +1}')$ in $T'$ containing only layer-2 vertices so that $r$ is the root of $T'$, $v = w_i'$ for some $1 \\leq i' \\leq {\\ell_{\\operatorname{pump}}} +1$, and the half-edge labels of $w_{ {\\ell_{\\operatorname{pump}}} +1}'$ have been fixed by $f$.\nSince $x \\in \\mathcal{V}' \\setminus \\mathcal{V}^\\ast$, we have $v \\neq w_{ {\\ell_{\\operatorname{pump}}} +1}'$. Hence $P = (v = w_i', w_{i+1}', \\ldots, w_{ {\\ell_{\\operatorname{pump}}} +1}')$ is a desired path with $2 \\leq k \\leq {\\ell_{\\operatorname{pump}}} + 1$ vertices.\n\n \n \n \n\n\\paragraph{Case 3} There is only one remaining neighbor of $v$ to consider. In view of the construction of $X_f(W_3^\\ast)$, there is a layer-3 vertex $v'$ adjacent to the vertex $r$, the root of the tree $T' \\in W_3^\\ast$ considered in the previous case. If the half-edge labels of $v'$ have been fixed by $f$, then $P = (v = w_i', w_{i-1}', \\ldots, w_{1}'=r, v')$ is a desired path.\nOtherwise, similar to the analysis in the previous case, we can find a path $P' = (v', \\ldots, w)$ connecting $v'$ to a vertex $w$ whose half-edge labels have been fixed by $f$. All vertices in $P'$ are of layer-3, and the number of vertices in $P'$ is within $[2, {\\ell_{\\operatorname{pump}}} + 1]$. \nCombining $P'$ with the path $(v = w_i', w_{i-1}', \\ldots, w_{1}'=r, v')$, we obtain the desired path $P$ whose number of vertices satisfies $2 \\leq k \\leq 2 {\\ell_{\\operatorname{pump}}} + 1$.\n \n \n\\end{proof}\n\n\n\\begin{figure}[h!]\n \\centering\n \\includegraphics[width= .75\\textwidth]{fig\/y2.pdf}\n \\caption{An illustration for the proof of \\cref{lem:find_path_1}.}\n \\label{fig:ell_full_construction1}\n\\end{figure}\n\nFor \\cref{lem:find_path_2}, note that if $\\mathtt{a}$ appears more than once in the multiset $x$, then we still have $\\mathtt{a} \\in x \\setminus \\{\\mathtt{a}\\}$.\n\n\n\\begin{lemma}\\label{lem:find_path_2}\nFor any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$, $\\mathtt{a}' \\in x' \\in \\mathcal{V}^\\ast$, and $0 \\leq t \\leq {\\ell_{\\operatorname{pump}}} - 1$, we have $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ for some $\\mathtt{b} \\in x \\setminus \\{\\mathtt{a}\\}$ and $\\mathtt{b}' \\in x' \\setminus \\{\\mathtt{a}'\\}$\n with $k = 2 {\\ell_{\\operatorname{pump}}} + 4 + t$.\n\\end{lemma}\n\\begin{proof}\nFor any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$, we can find a rooted tree $T_{x, \\mathtt{a}} \\in W_2^\\ast$ such that the path $(v_1$, $v_2$, $\\ldots$, $v_{ {\\ell_{\\operatorname{pump}}} +1})$ of the layer-1 vertices in $T_{x, \\mathtt{a}}$ where $v_1 = r$ is the root of $T_{x, \\mathtt{a}}$ satisfies the property that the half-edge labels of $v_{ {\\ell_{\\operatorname{pump}}} +1}$ have been fixed to $x$ by $f$ where the half-edge label for $\\{v_{ {\\ell_{\\operatorname{pump}}} }, v_{ {\\ell_{\\operatorname{pump}}} +1}\\}$ is in $x \\setminus \\{\\mathtt{a}\\}$.\n\nNow, observe that for any choice of $(\\Delta-2)( {\\ell_{\\operatorname{pump}}} +1)$ rooted trees \\[\\{T_{i,j}\\}_{1 \\leq i \\leq {\\ell_{\\operatorname{pump}}} +1, 1 \\leq j \\leq \\Delta-2}\\] in $W_2^\\ast$, there is a rooted tree $T' \\in W_3^\\ast = Z_f(W_2^\\ast)$ formed by appending $T_{i, 1}$, $T_{i,2}$, $\\ldots$, $T_{i, \\Delta-2}$ to $u_i$ in the path $(u_1$, $u_2$, $\\ldots$, $u_{ {\\ell_{\\operatorname{pump}}} +1})$ and fixing the half-edge labeling of $u_{ {\\ell_{\\operatorname{pump}}} +1}$ by $f$. All vertices in $(u_1, u_2, \\ldots, u_{ {\\ell_{\\operatorname{pump}}} +1})$ are of layer-2.\n\nWe choose any $T' \\in W_2^\\ast$ with $T_{i,j} = T_{x, \\mathtt{a}}$ and $T_{i',j'} = T_{x', \\mathtt{a}'}$ such that $i' - i = t + 1$.\nThe possible range of $t$ is $[0, {\\ell_{\\operatorname{pump}}} - 1]$.\nConsider any $T \\in W_L^\\ast$ that contains $T'$ as its subtree, and consider any correct labeling of $T$.\nThen the path $P$ resulting from concatenating the path $(v_{ {\\ell_{\\operatorname{pump}}} +1}, v_{ {\\ell_{\\operatorname{pump}}} }, \\ldots, v_1)$ in $T_{x, \\mathtt{a}}$, the path $(u_i, u_{i+1}, \\ldots, u_{i'})$, and the path $(v_1', v_2', \\ldots, v_{ {\\ell_{\\operatorname{pump}}} +1}')$ in $T_{x', \\mathtt{a}'}$\nshows that $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ \nfor $k = ( {\\ell_{\\operatorname{pump}}} +1) + (t+2) + ( {\\ell_{\\operatorname{pump}}} +1) = 2 {\\ell_{\\operatorname{pump}}} + 4 + t$. See \\cref{fig:ell_full_construction2} for an illustration.\n\nFor the rest of the proof, we show that the desired rooted tree $T_{x, \\mathtt{a}} \\in W_2^\\ast$ exists for any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$.\nFor any $x \\in \\mathcal{V}^\\ast$, we can find a bipolar tree\n\\[H_\\ast = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T_\\ast^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r) \\in Y_1^\\ast = Y_f(W_1^\\ast)\\]\nsuch that the vertex configuration of the root $r^m$ of $T_\\ast^m$ is fixed to be $x$ by the labeling function $f$.\nThen, for any $\\mathtt{a} \\in x$, at least one of the two rooted trees in $W_2^\\ast = Z_f(W_1^\\ast)$ resulting from cutting $H_\\ast$ satisfies the desired requirement. \n\\end{proof}\n\n\n\\begin{figure}[h!]\n \\centering\n \\includegraphics[width= .65\\textwidth]{fig\/y3.pdf}\n \\caption{An illustration for the proof of \\cref{lem:find_path_2}.}\n \\label{fig:ell_full_construction2}\n\\end{figure} \n\nIn the subsequent discussion, the length of a path refers to the number of edges in a path. In particular, the length of a $k$-vertex path is $k-1$.\n\nThe following lemma is proved by iteratively applying \\cref{lem:find_path_2} via intermediate vertex configurations $\\tilde{x} \\in \\mathcal{V}^\\ast$. \n\n\\begin{lemma}\\label{lem:find_path_3}\nFor any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$ and $\\mathtt{a}' \\in x' \\in \\mathcal{V}^\\ast$, we have $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ for some $\\mathtt{b} \\in x \\setminus \\{\\mathtt{a}\\}$ and $\\mathtt{b}' \\in x' \\setminus \\{\\mathtt{a}'\\}$\n for all $k \\geq 6 {\\ell_{\\operatorname{pump}}} + 10$.\n\\end{lemma}\n\\begin{proof}\nBy applying \\cref{lem:find_path_2} for three times,\nwe infer that \\cref{lem:find_path_2} also works for any path length $k - 1 = (k_1 - 1) + (k_2 - 1) + (k_3 - 1)$, where $k_i - 1 = 2 {\\ell_{\\operatorname{pump}}} + 3 + t_i$ for $0 \\leq t_i \\leq {\\ell_{\\operatorname{pump}}} -1$. \n\nSpecifically, we first apply \\cref{lem:find_path_2} to find a path realizing $(x,\\mathtt{b}) \\overset{k_1}{\\leftrightarrow} \\tilde{x}$ for some $\\tilde{x} \\in \\mathcal{V}^\\ast$, and for some $\\mathtt{b} \\in x \\setminus \\{\\mathtt{a}\\}$.\nLet $\\tilde{\\mathtt{a}} \\in \\tilde{x}$ be the real half-edge label used in the last vertex of the path. \nWe extend the length of this path by $k_2 - 1$ by attaching to it the path realizing $(\\tilde{x},\\tilde{\\mathtt{b}}) \\overset{k_2}{\\leftrightarrow} \\tilde{x}$, for some $\\tilde{\\mathtt{b}} \\in \\tilde{x} \\setminus \\{\\tilde{\\mathtt{a}}\\}$.\nFinally, let $\\tilde{\\mathtt{c}} \\in \\tilde{x}$ be the real half-edge label used in the last vertex of the current path. \nWe extend the length of the current path by $k_3$ by attaching to it the path realizing $(\\tilde{x},\\tilde{\\mathtt{d}}) \\overset{k_3}{\\leftrightarrow} (x', \\mathtt{b}')$, for some $\\tilde{\\mathtt{d}} \\in \\tilde{x} \\setminus \\{\\tilde{\\mathtt{x}}\\}$, and for some $\\mathtt{b}' \\in x' \\setminus \\{\\mathtt{a}'\\}$.\nThe resulting path realizes $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ for some $\\mathtt{b} \\in x \\setminus \\{\\mathtt{a}\\}$ and $\\mathtt{b}' \\in x' \\setminus \\{\\mathtt{a}'\\}$.\n\n\nTherefore, \\cref{lem:find_path_2} also works with path length of the form $k-1 = 6 {\\ell_{\\operatorname{pump}}} + 9 + t$, for any $t \\in [0, 3 {\\ell_{\\operatorname{pump}}} - 3]$. \n\nAny $k-1 \\geq 6 {\\ell_{\\operatorname{pump}}} + 9$ can be written as $k-1 = b(2 {\\ell_{\\operatorname{pump}}} +3) + (6 {\\ell_{\\operatorname{pump}}} + 9 + t)$, for some $t \\in [0, 3 {\\ell_{\\operatorname{pump}}} - 3]$ and some integer $b \\geq 0$.\nTherefore, similar to the above, we can find a path showing $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ by first applying \\cref{lem:find_path_2} with path length\n$2 {\\ell_{\\operatorname{pump}}} +3$ for $b$ times, and then applying the above variant of \\cref{lem:find_path_2} to extend the path length by $6 {\\ell_{\\operatorname{pump}}} + 9 + t$.\n\\end{proof}\n\nWe show that \\cref{lem:find_path_1,lem:find_path_3} imply that $\\mathcal{V}'$ is $\\ell$-full for some $\\ell$. \n\n\\begin{lemma}[$\\mathcal{V}'$ is $\\ell$-full] \\label{lem:find_path_4}\nThe set $\\mathcal{V}'$ is $\\ell$-full for $\\ell = 10 {\\ell_{\\operatorname{pump}}} + 10$. \n\\end{lemma}\n\\begin{proof}\nWe show that for any target path length $k-1 \\geq 10 {\\ell_{\\operatorname{pump}}} + 9$, and for any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$ and $\\mathtt{a}' \\in x' \\in \\mathcal{V}^\\ast$, we have $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} (x',\\mathtt{a}')$.\n\n\nBy \\cref{lem:find_path_1}, there exists a vertex configuration $\\tilde{x} \\in \\mathcal{V}^\\ast$ so that we can find a path $P$ realizing $(x,\\mathtt{a}) \\overset{\\ell_1}{\\leftrightarrow} \\tilde{x}$ for some path length $1 \\leq (\\ell_1-1) \\leq 2 {\\ell_{\\operatorname{pump}}} $.\nSimilarly, there exists a vertex configuration $\\tilde{x}' \\in \\mathcal{V}^\\ast$ so that we can find a path $P'$ realizing $(x',\\mathtt{a}') \\overset{\\ell_2}{\\leftrightarrow} \\tilde{x}'$ for some path length $1 \\leq (\\ell_2-1) \\leq 2 {\\ell_{\\operatorname{pump}}} $. \n\n\nLet $\\tilde{\\mathtt{a}}$ be the real half-edge label for $\\tilde{x}$ in $P$, and let $\\tilde{\\mathtt{a}}'$ be the real half-edge label for $\\tilde{x}'$ in $P'$.\nWe apply \\cref{lem:find_path_3} to find a path $\\tilde{P}$ realizing $(x,\\mathtt{b}) \\overset{\\tilde{\\ell}}{\\leftrightarrow} (x', \\mathtt{b}')$ for some $\\mathtt{b} \\in \\tilde{x} \\setminus \\{\\tilde{\\mathtt{a}}\\}$ and $\\mathtt{b}' \\in \\tilde{x}' \\setminus \\{\\tilde{\\mathtt{a}}'\\}$ with path length \n\\[\\tilde{\\ell}-1 = (k-1) - (\\ell_1-1) - (\\ell_2-1) \\geq 6 {\\ell_{\\operatorname{pump}}} + 9.\\] \nThe path formed by concatenating $P_1$, $\\tilde{P}$, and $P_2$ shows that $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} (x',\\mathtt{a}')$.\n\\end{proof}\n\n\n\n\n\n\n \n\n\n\n\\paragraph{A Feasible Labeling Function $f$ exists} We show that a feasible labeling function $f$ exists given that $\\Pi$ can be solved in $\\mathsf{LOCAL}(n^{o(1)})$ rounds. We will construct a labeling function $f$ in such a way each equivalence class in $\\textsf{Class}_1(W^\\ast)$ contains only rooted trees that admit legal labeling. That is, for each $c \\in \\textsf{Class}_1(W^\\ast)$, the mapping $h$ associated with $c$ satisfies $h(X) = \\textsf{YES}$ for some $X$.\n\n\n\nWe will consider a series of modifications in the construction \\[W \\rightarrow X_f(W) \\rightarrow Y_f(W) \\rightarrow Z_f(W)\\] that do not alter the sets of equivalence classes in these sets, conditioning on the assumption that all trees that we have processed admit correct labelings. That is, whether $f$ is feasible is invariant of the modifications.\n\n\n\\paragraph{Applying the Pumping Lemma} Let $w > 0$ be some target length for the pumping lemma. In the construction of $Y_f(W)$, when we process a bipolar tree \n\\[H = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r),\\] we apply \\cref{lem:pump} to the two subtrees $H^l = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l)$ and $H^r = (T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)$ to obtain two new bipolar trees $H_+^l$ and $H_+^r$.\nThe $s$-$t$ path in the new trees $H_+^l$ and $H_+^r$ contains $w + x$ vertices, for some $0 \\leq x < {\\ell_{\\operatorname{pump}}} $. The equivalence classes do not change, that is, $\\texttt{type}(H^l) = \\texttt{type}(H_+^l)$ and $\\texttt{type}(H^r) = \\texttt{type}(H_+^r)$.\n\n\nWe replace $H^l$ by $H_+^l$ and replace $H^r$ by $H_+^r$ in the bipolar tree $H$.\nRecall that the outcome of applying the labeling function $f$ to the root $r$ of the rooted tree $T_{ {\\ell_{\\operatorname{pump}}} +1}$ depends only on \\[\\texttt{type}(H^l), \\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m), \\texttt{type}(H^r),\\] so applying the pumping lemma to $H^l$ and $H^r$ during the construction of $Y_f(W)$ does not alter $\\textsf{Class}_1(T_\\ast^m)$ and $\\texttt{type}(H^\\ast)$ for the resulting bipolar tree $H^\\ast$ and its middle rooted tree $T_\\ast^m$, by \\cref{lem:labeling_function}.\n\n\n\n\\paragraph{Reusing Previous Trees} During the construction of the fixed point $W^\\ast$, we remember all bipolar trees to which we have applied the feasible function $f$.\n\nDuring the construction of $Y_f(W)$. Suppose that we are about to process a bipolar tree $H$, and there is already some other bipolar tree $\\tilde{H}$ to which we have applied $f$ before so that \n$\\texttt{type}(H^l)=\\texttt{type}(\\tilde{H}^l)$, $\\textsf{Class}_1(T_{i}^m) = \\textsf{Class}_1(\\tilde{T}_{i}^m)$ for each $1 \\leq i \\leq \\Delta-2$, and $\\texttt{type}(H^r)=\\texttt{type}(\\tilde{H}^r)$.\nThen we replace $H$ by $\\tilde{H}$, and then we process $\\tilde{H}$ instead.\n\nBy \\cref{lem:labeling_function}, this modification does not alter $\\textsf{Class}_1(T_\\ast^m)$ and $\\texttt{type}(H^\\ast)$ for the resulting bipolar tree $H^\\ast$.\n\n\\paragraph{Not Cutting Bipolar Trees} We consider the following different construction of $Z_f(W)$ from $Y_f(W)$. For each $H^\\ast \\in Y_f(W)$, we simply add two copies of $H^\\ast$ to $Z_f(W)$, one of them has $r = s$ and the other one has $r = t$.\nThat is, we do not cut the bipolar tree $H^\\ast$, as in the original construction of $Z_f(W)$.\n\nIn general, this modification might alter $\\textsf{Class}_1(Z_f(W))$.\nHowever, it is not hard to see that if all trees in $Y_f(W)$ already admit correct labelings, then $\\textsf{Class}_1(Z_f(W))$ does not alter after the modification.\n\nSpecifically, let $T$ be a rooted tree in $Z_f(W)$ with the above modification. \nThen $T$ is identical to a bipolar tree $H^\\ast=(T_1, T_2, \\ldots, T_k) \\in Y_f(W)$, where there exists some $1 < i < k$ such that the root $r_i$ of $T_i$ has its half-edge labels fixed by $f$.\nSuppose that the root of $T$ is the root $r_1$ of $T_1$.\nIf we do not have the above modification, then the rooted tree $T'$ added to $Z_f(W)$ corresponding to $T$ is $(T_1, T_2, \\ldots, T_i)$, where the root of $T'$ is also $r_1$. \n\nGiven the assumption that all bipolar trees in $Y_f(W)$ admit correct labelings, it is not hard to see that $\\textsf{Class}_1(T') = \\textsf{Class}_1(T)$. For any correct labeling $\\mathcal{L}'$ of $T'$, we can extend the labeling to a correct labeling $\\mathcal{L}$ of $T$ by labeling the remaining vertices according to any arbitrary correct labeling of $H^\\ast$. This is possible because the labeling of $r_i$ has been fixed. For any correct labeling $\\mathcal{L}$ of $T$, we can obtain a correct labeling $\\mathcal{L}'$ of $T'$ by restricting $\\mathcal{L}$ to $T'$.\n\n\n\n\n\n\n\\paragraph{Simulating a $\\mathsf{LOCAL}$ Algorithm} We will show that there is a labeling function $f$ that makes all the rooted trees in $W^\\ast$ to admit correct solutions, where we apply the above three modifications in the construction of $W^\\ast$. In view of the above discussion, such a function $f$ is feasible.\n\nThe construction of $W^\\ast$ involves only a finite number $k$ of iterations of $Z_f$ applied to some previously constructed rooted trees. If we view the target length $w$ of the pumping lemma as a variable, then the size of a tree in $W^\\ast$ can be upper bounded by $O(w^{k})$. \n\n\nSuppose that we are given an arbitrary $n^{o(1)}$-round $\\mathsf{LOCAL}$ algorithm $\\mathcal{A}$ that solves $\\Pi$. It is known~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup,chang2020n1k_speedups} that randomness does not help for LCL problems on bounded-degree trees with round complexity $\\Omega(\\log n)$, so we assume that \n$\\mathcal{A}$ is deterministic.\n\nThe runtime of $\\mathcal{A}$ on a tree of size \n$O(w^{k})$ can be made to be at most $t = w\/10$ if $w$ is chosen to be sufficiently large. \n\nWe pick the labeling function $f$ as follows. \nWhenever we encounter a new bipolar tree $H$ to which we need to apply $f$, we simulate $\\mathcal{A}$ locally at the root $r^m$ of the middle rooted tree $T^m$ in $H$. Here we assume that the pumping lemma was applied before simulating $\\mathcal{A}$. Moreover, when we do the simulation, we assume that the number of vertices is $n = O(w^{k})$. \n\nTo make the simulation possible, we just locally generate arbitrary distinct identifiers in the radius-$t$ neighborhood $S$ of $r^m$. This is possible because $t = w\/10$ is so small that for any vertex $v$ in any tree $T$ constructed by recursively applying $Z_f$ for at most $k$ iterations, the radius-$(t+1)$ neighborhood of $v$ intersects at most one such $S$-set. Therefore, it is possible to complete the identifier assignment to the rest of the vertices in $T$ so that for any edge $\\{u,v\\}$, the set of identifiers seen by $u$ and $v$ are distinct in an execution of a $t$-round algorithm. \n\nThus, by the correctness of $\\mathcal{A}$, the labeling of all edges $\\{u,v\\}$ are configurations in $\\mathcal{E}$ and the labeling of all vertices $v$ are configurations in $\\mathcal{V}$.\nOur choice of $f$ implies that all trees in $W^\\ast$ admit correct labeling, and so $f$ is feasible. Combined with \\cref{lem:find_path_4}, we have proved that $\\mathsf{LOCAL}(n^{o(1)})$ $\\Rightarrow$ $\\ell$-full.\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Concluding Remarks}\\label{sec:concludingRemarksTrees}\n\nHere we collect some open problems we find interesting. \n\n\n\\begin{enumerate}\n \\item What is the relation of $\\mathsf{ffiid}, \\mathsf{MEASURE}, \\mathsf{fiid}$? In particular is perfect matching in the class $\\mathsf{ffiid}$ or $\\mathsf{MEASURE}$? \n \\item Are there LCL problems whose uniform complexity is nontrivial and slower than $\\operatorname{poly}\\log 1\/\\varepsilon$?\n \n \n \n \n \n \n \\item Suppose that $\\Pi_G$ is a homomorphism LCL, is it true that $\\Pi_G\\in\\mathsf{BOREL}$ if and only if $\\Pi_G\\in O(\\log^*(n))$?\n \\item Are there finite graphs of chromatic number bigger than $2\\Delta-2$ such that the corresponding homomorphism LCL is not in $\\mathsf{BOREL}$? See also Remark \\ref{r:hedet}.\n\\end{enumerate}\n\n\\subsection*{Acknowledgements} \nWe would like to thank Anton Bernshteyn, Endre Cs\u00f3ka, Mohsen Ghaffari, Jan Hladk\u00fd, Steve Jackson, Alexander Kechris, Edward Krohne, Oleg Pikhurko, Brandon Seward, Jukka Suomela, and Yufan Zheng for insightful discussions. \nWe especially thank Yufan Zheng for suggesting potentially hard instances of homomorphism LCL problems.\n\n\n\\bibliographystyle{alpha}\n\n\\section{Introduction}\n\nIn this paper, we study local problems on grid graphs from three different perspectives: (a) the perspective of the theory of distributed algorithms, (b) the perspective of the theory of random processes, (c) the perspective of the field of descriptive combinatorics. \nAlthough different communities in these three fields ask very similar questions, no systematic connections between them were known, until the insightful work of Bernshteyn \\cite{Bernshteyn2021LLL} who applied results from the field (a) to the field (c). In this work, our most important results follow by applying techniques from the field (a) to (b), but perhaps more importantly, we present all three different perspectives as a part of one common theory (see \\cref{fig:big_picture_grids}). \n\nIn this introductory section, we first quickly present the three fields and define the object of interest, local problems, a class of problems that contains basic problems like vertex or edge coloring, matchings, and others. \nThen, we briefly present our results, most of which can be seen as inclusions between different classes of local problems shown in \\cref{fig:big_picture_grids}. \n\n\\paragraph{Distributed Computing}\n\nThe $\\mathsf{LOCAL}$ model of computing \\cite{linial1987LOCAL} was motivated by understanding distributed algorithms in huge networks. \nAs an example, consider the network of all wifi routers, where two routers are connected if they are close enough to exchange messages. In this case, they should better communicate with user devices on different channels, so as to avoid interference. In the graph-theoretic language, we want to color the network. Even if the maximum degree is $\\Delta$ and we want to color it with $\\Delta+1$ colors, the problem remains nontrivial, because the decision of each node should be done after few communication rounds with its neighbors. \n\nThe $\\mathsf{LOCAL}$ model formalizes this setup: we have a large network, with each node knowing its size, $n$, and perhaps some other parameter like the maximum degree $\\Delta$. In case of randomized algorithms, each node has an access to a random string, while in case of deterministic algorithms, each node starts with a unique identifier from a range of size polynomial in $n$. \nIn one round, each node can exchange any message with its neighbors and can perform an arbitrary computation. We want to find a solution to a problem in as few communication rounds as possible. \nImportantly, there is an equivalent view of $t$-rounds $\\mathsf{LOCAL}$ algorithms: such an algorithm is simply a function that maps $t$-hop neighbourhoods to the final output. An algorithm is correct if and only if applying this function to each node solves the problem. \n\nThere is a rich theory of distributed algorithms and the local complexity of many problems is understood. One example: if the input graph is the $d$-dimensional torus (think of the grid $\\mathbb{Z}^d$ and identify nodes with the same remainder after dividing the coordinate with $n^{1\/d}$), a simple picture emerges. Any local problem (we define those formally later) is solvable with local complexity $O(1)$, $\\Theta(\\log^* n)$, or we need to see the whole graph to solve it, i.e., the local complexity is $\\Theta(n^{1\/d})$ (see \\cref{thm:classification_of_local_problems_on_grids}). For example, a proper coloring of nodes with $4$ colors has local complexity $\\Theta(\\log^* n)$, while coloring with $3$ colors has already complexity $\\Theta(n^{1\/d})$. \n\n\\paragraph{Finitary factors of iid processes}\n\\todo{zeptat se Spinky}\nConsider the following example of the so-called anti-ferromagnetic Ising model. \nImagine atoms arranged in a grid $\\mathbb{Z}^3$. Each atom can be in one of two states and there is a preference between neighboring atoms to be in a different state. We can formalize this by defining a potential function $V(x,y)$ indicating whether $x$ and $y$ are equal. \nThe Ising model models the behavior of atoms in the grid as follows: for the potential $V$ and inverse temperature $\\beta$, if we look at the atoms, we find them in a state $\\Lambda$ with probability proportional to $\\exp(-\\beta \\sum_{u \\sim v \\in E(\\mathbb{Z}^d)} V(\\Lambda(u), \\Lambda(v)))$. If $\\beta$ is close to zero, we sample essentially from the uniform distribution. If $\\beta$ is infinite, we sample one of the two $2$-colorings of the grid. \n\nTo understand better the behavior of phase transitions, one can ask the following question. For what values of $\\beta$ can we think of the distribution $\\mu$ over states from the Ising model as follows: there is some underlying random variable $X(u)$ at every node $u$, with these variables being independent and identically distributed. We can think of $\\mu$ as being generated by a procedure, where each node explores the random variables in its neighborhood until it decides on its state, which happens with probability $1$. If this view of the underlying distribution is possible, we say it is a finitary factor of an iid process (ffiid). \nIntuitively, being an ffiid justifies the intuition that there is a correlation decay between different sites of the grid. To understand ffiids, we try to analyze solutions to which local problems can be sampled from an ffiid distribution (2-coloring of the grid is not among them). \n\nAlthough this definition was raised in a different context then distributed computing, it is actually almost identical to the definition of a distributed algorithm! The differences are two. First, the ffiid procedure does not know the value of $n$ since, in fact, it makes most sense to consider it in infinite graphs (like the infinite grid $\\mathbb{Z}^d$). Second, the natural measure of the complexity is different: it is the speed of the decay of the random variable corresponding to the time until the local algorithm finds the solution. \nNext, note that the algorithm has to return a correct solution with probability $1$, while the usual randomized local algorithm can err with a small probability. We call ffiids \\emph{oblivious local algorithms} to keep the language of distributed computing (see \\cref{sec:preliminaries}).\nWe define the class $\\mathsf{TAIL}(f(r))$ as the class of problems solvable by an oblivious local algorithm if the probability of the needed radius being bigger than $r$ is bounded by $1\/f(r)$ (see \\cref{fig:big_picture_grids}). \n\n\\paragraph{Descriptive Combinatorics}\n\nA famous Banach-Tarski theorem states that one can split a three-dimensional unit volume ball into five pieces, translate and rotate each piece and end up with two unit volume balls. \nThis theorem serves as an introductory warning in the beginner measure theory courses: not every subset of $\\mathbb{R}^3$ can be measurable.\nSimilarly surprising was when in 1990 Laczkovich answered the famous Circle Squaring Problem of Tarski\\todo{cite laczkovich}: A square of unit area can be decomposed into finitely many pieces such that one can translate each of them to get a circle of unit area.\nUnlike in Banach-Tarski paradox, this result was improved in recent years to make the pieces measurable in various sense \\cite{OlegCircle, Circle, JordanCircle}.\n\n\nThis theorem and its subsequent strengthenings are the highlights of a field nowadays called descriptive combinatorics \\cite{kechris_marks2016descriptive_comb_survey,pikhurko2021descriptive_comb_survey} that has close connections to distributed computing as was shown in an insightful paper by Bernshteyn \\cite{Bernshteyn2021LLL}. \nGoing back to the Circle Squaring Problem, the proofs of its versions first choose a clever set of constantly many translations $\\tau_1, \\dots, \\tau_k$. Then, one constructs a bipartite graph between the vertices of the square and the cycle and connects two vertices $x, y$ if there is a translation $\\tau_i$ mapping the point $x$ in the square to the point $y$ in the cycle. The problem can now be formulated as a perfect matching problem and this is how combinatorics enters the picture. In fact, versions of a matching or flow problems are then solved on the grid graph $\\mathbb{Z}^d$ using a trick called \\emph{toast construction} also used in the area of the study of ffiids \\cite{HolroydSchrammWilson2017FinitaryColoring,ArankaFeriLaci,Spinka}. \n\nWe discuss the descriptive combinatorics classes in a more detail in the expository \\cref{sec:OtherClasses} but we will not need them in the main body. Let us now give an intuitive definition of the class $\\mathsf{BOREL}$. \nBy a result of Chang and Pettie \\cite{chang2016exp_separation}, we know that the maximal independent set problem is complete for the class $\\mathsf{LOCAL}(O(\\log^* n))$ in the sense that this problem can be solved in $O(\\log^* n)$ rounds and any problem in that class can be solved in constant number of rounds if we are given access to an oracle computing MIS (we make this precise in \\cref{rem:chang_pettie_decomposition}). Imagine now that the input graph is infinite and our local algorithm has access to MIS oracle and local computations but now both can be used countably many times instead of just once or, equivalently, fixed number of times as is the case with the class $\\mathsf{LOCAL}(O(\\log^* n))$. That is, we can think of an algorithm that constructs MIS of larger and larger powers of the input graph and uses this infinite sequence to solve an LCL. Those constructions are in the class $\\mathsf{BOREL}$ from \\cref{fig:big_picture_grids}. In fact, the exact relationship between this definition and full power of $\\mathsf{BOREL}$ is not clear.\nHowever, it should give an idea about the difference between descriptive combinatorics and distributed computing -- from the descriptive combinatorics perspective, reasonable classes are closed under repeating basic operations countably many times.\n\n\n\n\n\\paragraph{Locally Checkable Problems on Grids as a ``Rosetta Stone''}\n\nLocally checkable problems (LCLs) are graph problems where a correctness of a solution can be checked locally. That is, there is a set of local constraints such that if they are satisfied at every node, the solution to a problem is correct. \nAn examples of such constraints include vertex coloring, edge coloring, or perfect matching. We also allow inputs on each node, so the list coloring is also an example of LCL (given that the total number of colors in all lists is finite). One should think of this class as an analogue of the class NP from complexity theory. \n\nIn this paper, we study locally checkable problems on grid graphs in the three settings introduced above. This approach reveals number of connections. For example, in \\cref{fig:big_picture_grids}, we can see that the two classes of local complexities $\\mathsf{LOCAL}(O(1))$ and $\\mathsf{LOCAL}(O(\\log^* n))$ are equivalent to different natural classes from the perspective of factors of iid and descriptive combinatorics. On the other hand, those two worlds are richer. Many of our discussions in fact apply to graphs of bounded growth, but the grids serve as the simplest, but still quite nontrivial model in which one can try to understand all three settings together (the case of lines is understood quite well \\cite{grebik_rozhon2021LCL_on_paths} and we study aspects of the connection for trees in a parallel paper \\cite{brandt_chang_grebik_grunau_rozhon_vidnyaszky2021LCLs_on_trees_descriptive}). \n\n\\begin{landscape}\n\\begin{figure}\n \\centering\n \\includegraphics[width = 1.5\\textwidth]{fig\/big_picture_grids.png}\n \\caption{The big picture of local complexity classes on grids from the perspective of distributed algorithms, factors of iid, and descriptive combinatorics. Canonical example of each class is on the left side. }\n \\label{fig:big_picture_grids}\n\\end{figure}\n\\end{landscape}\n\n\n\\subsection{Our Contribution }\n\nOur contribution can be mostly summed up as several new arrows in \\cref{fig:big_picture_grids}. In general, this clarifies several basic questions regarding the connections of the three fields and helps to view all three fields as part of a general theory of locality. \n \n\n\nFor example, consider the following question from the paper of Brandt et al. \\cite{brandt_grids} that considered LCL problems on grids: \n\\begin{quote}\n\\emph{\nHowever, despite the fact that techniques seem to translate between finitary colorings and distributed complexity, it remains unclear how to directly translate results from one setting to the other in a black-box manner; for instance, can we derive the lower bound for 3-colouring from\nthe results of Holroyd et al., and does our complexity classification imply answers to the open questions they pose?\n}\n\\end{quote}\nWe show how to 'directly translate' results between the setting of distributed algorithms and finitary factors of iid in \\cref{subsec:preliminaries_tail,subsec:speedup_warmup}. This indeed answers two open questions of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n\nTo understand the connection between the world of factors of iids and descriptive combinatorics, we formalize the notion of toast construction used in both setups. We prove general results connecting the toast constructions with oblivious local algorithms and this enables us to understand the $\\mathsf{TAIL}$ complexity of a number of concrete problems like vertex $3$-coloring, edge $4$-coloring, perfect matchings and others. \n\nIn general, we hope that the concept of $\\mathsf{TAIL}$ complexities will find other usecases. We discuss in \\cref{subsec:congestLCA} how it relates to studied notions of oblivious algorithms, energy complexity, or its relevance to the $\\mathsf{LCA}$ model.\\todo{Oleg se ptal jestli uz ma tady vedet co je LCA a jestli to nekde neprehlid}\nIn fact, the $\\mathsf{TAIL}$ complexity approach was already considered e.g. in the celebrated MIS paper of \\cite{ghaffari2016MIS} (see the introduction of his paper) and the $\\mathsf{TAIL}$ bound on the MIS algorithm there may be the reason of its wide applicability. \n\nWe now list concrete results that we prove (or they follow from the work of others, see \\cref{fig:big_picture_grids}). \n\n\\paragraph{Classification Theorems for $\\mathsf{TAIL}$}\n\nWe prove several results that one can view as arrows in \\cref{fig:big_picture_grids}. Those can be viewed as an analogy to the line of work aiming for the classification of $\\mathsf{LOCAL}$ complexities discussed in \\cref{subsec:preliminaries_local}. \n\n\\begin{itemize}\n \\item $\\mathsf{TAIL}(\\mathrm{tower}(\\omega(r))) = \\mathsf{LOCAL}(O(1))$\\hspace*{1cm} We show that any oblivious local algorithm with tail decay faster than tower function can be in fact solved in $O(1)$ local complexity. This answers a question of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n \\item $\\mathsf{LOCAL}(O(\\log^* n)) = \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r))) = \\mathsf{CONTINUOUS}$\\hspace*{1cm} There are three natural definitions of classes of basic symmetry breaking problems in all three considered fields and they are the same \\cite{GJKS,HolroydSchrammWilson2017FinitaryColoring}. \n As a direct application one gets that from the three proofs that $3$-coloring of grids is not in $\\mathsf{LOCAL}(O(\\log^* n))$ \\cite{brandt_grids}, not in $ \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$ \\cite{HolroydSchrammWilson2017FinitaryColoring} and not in $ \\mathsf{CONTINUOUS}$ \\cite{GJKS} any one implies the others. These three classes are the same in a much more general setup (see \\cref{subsec:CONT}). \n \\item $\\mathsf{TAIL}(\\omega(r^{d-1})) = \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$\\hspace*{1cm} We prove that any problem that can be solved with an oblivious local algorithm with tail decay faster than $r^{d-1}$, can in fact be solved with a tower tail decay. This answers a question of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n \\item Similarly to the time hierarchy theorems of the $\\mathsf{LOCAL}$ model \\cite{chang2017time_hierarchy}, we study the range of possible $\\mathsf{TAIL}$ complexities and give examples of problems LCL problems with $\\mathsf{TAIL}$-complexities $\\Theta(r^{j - o(1)})$ for any $1 \\le j \\le d-1$. \n\\end{itemize}\n\n\\paragraph{$\\TOAST$ constructions}\nWe formally define $\\TOAST$ algorithms (we mentioned earlier that these constructions are used in descriptive combinatorics and the theory of random processes) and prove results about them. \nLet us first explain the toast constructions in a more detail. A $q$-toast in the infinite graph $\\mathbb{Z}^d$ is a collection $\\mathcal{D}$ of finite connected subsets of $\\mathbb{Z}^d$ that we call \\emph{pieces} such that (a) for any two nodes of the graph we can find a piece that contains both,\\todo{it is enough that every vertex is contained in a piece} (b) boundaries of any two different pieces are at least distance $q$ apart. \nThat is, imagine a collection of bigger and bigger pieces that form a laminar family (only inclusions, not intersections) that, in the limit, covers the whole infinite graph.\n\nThis structure can be built in the descriptive combinatorics setups \\cite{FirstToast,GJKS,Circle} but also with oblivious local algorithms \\cite{HolroydSchrammWilson2017FinitaryColoring,Spinka}. In fact, we show that one can construct the structure with $1\/r^{1-o(1)}$ $\\mathsf{TAIL}$ decay in \\cref{thm:toast_in_tail}. This is best possible up to the $o(1)$ loss in the exponent. \nIf one has an inductive procedure building the toast, we can try to solve an LCL problem during the construction. That is, whenever a new piece of the toast is constructed, we define how the solution of the problem looks like at this piece. This is what we call a $\\TOAST$ algorithm. \nThere is a randomized variant of the $\\TOAST$ algorithm and a lazy variant. In the lazy variant one can postpone the decision to color given piece, but it needs to be colored at some point in the future. \nNext, we prove several relations between the $\\TOAST$ classes (see \\cref{fig:big_picture_grids}). \n\n\\begin{itemize}\n \\item $\\TOAST \\subseteq \\mathsf{TAIL}(r^{1-o(1)})$\\hspace*{1cm} This result allows us to translate toast-algorithm constructions used in descriptive combinatorics to the $\\mathsf{TAIL}$-setting. This includes problems like vertex $3$-coloring (this answers a question of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}), edge $2d$-coloring, or perfect matching. \n \\item $\\RTOAST = \\TOAST$\\hspace*{1cm} This derandomization result for example implies on the first glance surprising inclusion $\\mathsf{FINDEP} \\subseteq \\mathsf{BOREL}$ -- these two classes are discussed in \\cref{sec:OtherClasses}. \n \\item $\\mathsf{TAIL} = \\RLTOAST$\\hspace*{1cm} Here, $\\mathsf{TAIL}$ is defined as the union of $\\mathsf{TAIL}(f)$ over all functions $f \\in \\omega(1)$. This result shows how the language of toast constructions can in fact characterize the whole hierarchy of constructions by oblivious local algorithms. \n \\item We show that the problem of $2d$ edge coloring of $\\mathbb{Z}^d$ is in $\\TOAST$, hence in $\\mathsf{BOREL}$ which answers the question of Gao et al \\cite{GJKS}. Independently, this was proven by \\cite{ArankaFeriLaci}. \n\\end{itemize}\n\n\\subsection{Roadmap}\n\nIn \\cref{sec:preliminaries} we recall the basic definitions regarding the $\\mathsf{LOCAL}$ model and explain the formal connection between finitary factors of iids. \nNext, in \\cref{sec:speedup_theorems} we prove the basic speedup theorems for $\\mathsf{TAIL}$. \nIn \\cref{sec:Toast} we define $\\TOAST$ algorithms and relate it to the $\\mathsf{TAIL}$ model. \nIn \\cref{sec:SpecificProblems} we discuss what our results mean for specific local problems and in \\cref{sec:Inter} we construct problems with $r^{j-o(1)}$ $\\mathsf{TAIL}$ complexities for $1 \\le j \\le d-1$. Finally, in \\cref{sec:OtherClasses} we discuss connection of our work to other classes of problems and finish with open questions in \\cref{sec:open_problems}. \n\n\n\n\n\n\\begin{comment}\nHonza - Introduction\n\nA \\emph{locally checkable problems (LCL) $\\Pi$ on a family of graphs $\\mathbf{G}$} is a procedure that checks whether a given ``candidate'' vertex coloring of $G\\in \\mathbf{G}$ satisfies a local constraints given by $\\Pi$\\footnote{\nFormally, $\\Pi$ is a triplet $\\Pi=(b,t,\\mathcal{P})$, where $b$ is a set of colors, $t\\in \\mathbb{N}$ is the locality of $\\Pi$ and $\\mathcal{P}$ is a set of allowed configurations around each vertex of locality at most $t$.}.\nMore concretely, given a graph $G=(V(G),E(G))\\in \\mathbf{G}$ and a ``candidate'' coloring $c:V(G)\\to b$, we say that $c$ is a \\emph{$\\Pi$-coloring} if for every $v\\in V(G)$ the restriction of $c$ to $t$-neighborhood of $v$ is in $\\mathcal{P}$, where $t$ is the locality of $\\Pi$ and $\\mathcal{P}$ is the set of constraints of $\\Pi$.\n\nIn this paper the family $\\mathbf{G}$ always refers to graphs that look locally like the standard Cayley graph of $\\mathbb{Z}^d$, $\\mathbb{Z}^d$, for some $d\\ge 1$.\nThis includes $\\mathbb{Z}^d$ as well as various finite $d$-dimensional tori with oriented edges that are labeled by the standard generators of $\\mathbb{Z}^d$.\nThe dimension $d$ is always assumed to be bigger than $1$, since for $d=1$ there is a full classification of LCL problems.\nWe omit mentioning ${\\bf G}$ and refer to all LCL problems on the family ${\\bf G}$ simply as \\emph{LCL problems (on $\n\\mathbb{Z}^d$)}.\nWhen we want to stress that a given LCL problem is on $\\mathbb{Z}^d$ we write $d$ in the superscript, i.e., $\\Pi^d$.\n\nWe study a complexity of LCL problems in various settings.\nIn the most abstract terms, a complexity of a given LCL problem $\\Pi$ is measured in terms of how difficult it is to \\emph{transform} a \\emph{given labeling} of a graph into a $\\Pi$-coloring.\nThe different contexts varies in what is considered to be an allowed starting labeling and what is allowed transformation.\nMost relevant for our investigation is the following probability notion.\nStart with a graph $G$ and assign to each vertex $v\\in V(G)$ a real number from $[0,1]$ uniformly and independently, i.e., \\emph{iid labeling} of $G$.\nA function $F$ that transforms these iid labels in a measurable and equivaraint way into a different coloring is called \\emph{factor of iid (fiid)}.\nHere equivariant refers to the group of all symmetries of $G$.\nIt is not hard to see that the group of all symmetries of $\\mathbb{Z}^d$ is the group $\\mathbb{Z}^d$ itself that acts by shift-translations.\nA fiid $F$ on $\\mathbb{Z}^d$ is called \\emph{finitary fiid (ffiid)} if for almost every iid labeling of $\\mathbb{Z}^d$ the color $F(x)({\\bf 0}_d)$ only depends on the labels of $x$ on some finite neighborhood of ${\\bf 0}_d$.\nThat is to say, if $x$ and $x'$ agrees on this neighborhood, then $F(x)({\\bf 0}_d)=F(x')({\\bf 0}_d)$.\nGiven ffiid $F$, we write $R_\\mathcal{A}$ for the random variable that codes minimum diameter of such neighborhood and we refer to it as the coding radius of $F$.\n\nWe say that a ffiid $F$ \\emph{solves an LCL problem $\\Pi$} if it produces a $\\Pi$-coloring for almost every iid labeling of $\\mathbb{Z}^d$.\nan LCL problem $\\Pi$ is in the class $\\mathsf{TAIL}$ if $\\Pi$ admits a ffiid solution.\nFollowing Holroyd, Schramm and Wilson \\cite{HolroydSchrammWilson2017FinitaryColoring}, the complexity of LCL problems in $\\mathsf{TAIL}$ can be measured on a finer scale by the properties of the coding radius $R_\\mathcal{A}$, where $F$ runs over all ffiid's that solve the given problem.\nSimilarly, we define when fiid $F$ solves $\\Pi$ and the class $\\mathsf{fiid}$.\n\nBefore we move to distributed algorithm we recall what is known about the notorious class of LCL problems that are the \\emph{proper vertex $k$-colorings of $\\mathbb{Z}^d$}, $\\Pi^d_{\\chi,k}$, where $k>1$, in the context of ffiid.\nIt is known that if $d>1$ and $k>3$, then $\\Pi^d_{\\chi,k}$ admits ffiid solution with a tower tail decay, see Holroyd et.al~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\nIf $d>0$ and $k=2$, then $\\Pi^d_{\\chi,2}$ is not even in $\\mathsf{fiid}$.\nThe most interesting case is when $d>1$ and $k=3$.\nIt was shown by Holroyd et al~\\cite{HolroydSchrammWilson2017FinitaryColoring} that in that case $\\Pi^d_{\\chi,3}$ is in the class $\\mathsf{TAIL}$ and every ffiid that solves $\\Pi^d_{\\chi,3}$ must have infinite the $2$-moment of its coding radius. \nMoreover, they found a ffiid that solve $\\Pi^d_{\\chi,3}$ and has finite some moments of its coding radius.\nHere we strengthen their result by showing that ${\\bf M}(\\Pi^d_{\\chi,3})=1$ and every ffiid solution of $\\Pi^d_{\\chi,3}$ has infinite the $1$-moment of its coding radius.\n\nIn the context of deterministic distributed computing the setting is the following.\nStart with finite graph $G$ on $n$-vertices (in our case $n$ is big enough and $G$ is a torus) and assign to every vertex of $G$ its unique identifier from $[n]$.\\footnote{In generality, the identifier could be assigned from some domain that has size polynomial in $n$.}\nA \\emph{local algorithm $\\mathcal{A}_n$ of locality $t_n$} is a rule that computes color of a vertex $v$ based only on its $t$-neighborhood and the unique identifiers in this neighborhood.\nWe say that $\\mathcal{A}_n$ \\emph{solves} $\\Pi$ if for every (admissible)\\footnote{As a torus $G$ needs to be wider than the locality of $\\mathcal{A}_n$} graph $G$ on $n$-vertices and any assignment of unique identifiers to $G$ the rule $\\mathcal{A}_n$ produces a $\\Pi$-coloring.\nWe say that $\\Pi$ is in the class $\\mathsf{LOCAL}(O(\\log^*(n)))$ if there is a sequence of local rules \\todo{co je to local rule?}$\\{\\mathcal{A}_n,t_n\\}_{n\\in \\mathbb{N}}$ such that $\\mathcal{A}_n$ has locality $t_n$, solves $\\Pi$ for every $n\\in\\mathbb{N}$ and $n\\mapsto t_n$ is in the class $O(\\log^*(n))$.\nThe reason why we picked $\\log^*(n)$ is that there is a complete classification of LCL problems on oriented labeled grids.\nRoughly speaking, an LCL problem $\\Pi$ is either \\emph{trivial}, i.e., there is a single color that produces $\\Pi$-coloring, \\emph{local}, i.e., there is an algorithm that needs locality $\\Theta(\\log^*(n))$, or \\emph{global}, i.e., every local algorithm that solves $\\Pi$ needs to see essentially the whole graph.\\footnote{Same classification holds in the context of randomized distributed algorithms.}\n\nWe use this classification to provide several speed-up results in the context of ffiid solutions.\nWe note that the connection between distributed computing and probability setting was previously observed by several authors.\nFor example, it follows easily from the result of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring} that local coloring problems that are not global in the distributed setting admit ffiid solution with tower tail decay.\nIn fact, Brandt et al. \\cite{brandt_grids} asked:\n\\begin{quote}\n\\emph{\nHowever, despite the fact that techniques seem to translate between finitary colourings\nand distributed complexity, it remains unclear how to directly translate results from one setting to\nthe other in a black-box manner; for instance, can we derive the lower bound for 3-colouring from\nthe results of Holroyd et al. [24], and does our complexity classification imply answers to the open\nquestions they pose?\n}\n\\end{quote}\nAs one of our main result we address their question.\n\nThe most abstract setting that we consider is Borel combinatorics.\nIn this context one studies graphs that additionally posses a completely metrizable and separable topological structure (or a standard Borel structure) on its vertex set.\nIt follows from the work of Seward and Tucker-Drob \\cite{ST} and Gao, Jackson, Krohne and Seward \\cite{GJKS} that the canonical object to investigate is the graph structure on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$, the free part of the shift action $\\mathbb{Z}^d\\curvearrowright 2^{\\mathbb{Z}^d}$, induced by the standard generating set of $\\mathbb{Z}^d$. \nThe topological and Borel structure on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ is the one inherited from the product topology on $2^{\\mathbb{Z}^d}$.\nWe say that an LCL problem $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ if there is a continuous and equivariant map from $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ into the space of $\\Pi$-colorings.\nThe class $\\mathsf{BOREL}$ is defined similarly, only the requirement on the map to be continuous is relaxed to be merely Borel measurable.\nBernshteyn \\cite{Bernshteyn2021LLL} proved that LCL problems that are not global in the distributed setting admit continuous and therefore Borel solution.\nIn another words, he showed that the class $\\mathsf{LOCAL}(O(\\log^*(n)))$ is contained in $\\mathsf{CONTINUOUS}$.\nIt is not hard to show that the other inclusion follows from the Twelve Tile Theorem of Gao et al \\cite{GJKS}, i.e., the class $\\mathsf{CONTINUOUS}$ in fact coincide with the class $\\mathsf{LOCAL}(O(\\log^*(n)))$.\nHere we investigate in greater detail what is the connection between existence of a Borel solution and a tail behavior of ffiid.\nIn fact all known Borel constructions uses the following structure that can be rephrased without mentioning topology or Borel measurability.\n\nA \\emph{$q$-toast on $\\mathbb{Z}^d$} is a collection $\\mathcal{D}$ of finite connected subsets of $\\mathbb{Z}^d$ such that\n\\begin{itemize}\n \\item for every $g,h\\in \\mathbb{Z}^d$ there is $D\\in\\mathcal{D}$ such that $g,h\\in D$,\n \\item if $D,E\\in \\mathcal{D}$ and $d(\\partial D,\\partial E)\\le q$, then $D=E$.\n\\end{itemize}\nWe say that an LCL problem $\\Pi$ is in the class $\\TOAST$ if there is $q\\in \\mathbb{N}$ and an equivariant rule $\\Omega$ that extends partial colorings such that if we apply $\\Omega$ on every $q$-toast $\\mathcal{D}$ inductively on the depth, then we produce a $\\Pi$-coloring.\nIt is not known whether the classes $\\mathsf{BOREL}$ and $\\TOAST$ coincide, however, it follows from the work of Gao et al.~\\cite{GJKS2} that $\\TOAST\\subseteq \\mathsf{BOREL}$.\nThis is proved by constructing a toast structure in a Borel way on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$, see Marks and Unger~\\cite{Circle} for application and proof of this fact.\nThe specific problems that belong to the class $\\TOAST$ include $3$-vertex coloring, perfect matching, $2d$-edge coloring, Schreier decoration (for non-oriented grids), some rectangular tilings, space filling curve etc.\nThese results were proved by many people in slightly different language, \\cite{GJKS2, ArankaFeriLaci, spencer_personal}.\nWe utilize here all the constructions and show that these all belong to $\\TOAST$.\n\nWe connect the classes $\\TOAST$ and $\\mathsf{TAIL}$ as follows.\nWe show that for every $q\\in \\mathbb{N}$ and $\\epsilon>0$ there is a ffiid $F$ that produces $q$-toast that has the $(1-\\epsilon)$-moment finite.\nA slightly stronger result implies that, in fact, ${\\bf M}(\\Pi)\\ge 1$, whenever an LCL problem $\\Pi$ is in the class $\\TOAST$.\nThis shows that ${\\bf M}(\\Pi)\\ge 1$ for all the problems mentioned above.\nWe note that in the special case when $d=2$ and the corresponding LCL problem $\\Pi$ does not admit ffiid solution with a tower tail decay, we deduce that ${\\bf M}(\\Pi)=1$.\nThis follows from our speed-up result in the context of ffiid.\n\nTraditionally there are two more classes that are investigated in Borel combinatorics, $\\mathsf{BAIRE}$ and $\\mathsf{MEASURE}$.\nTo define the class $\\mathsf{BAIRE}$ we replace the topological requirement on continuity by \\emph{Baire measurability}, i.e., we require a solution to be continuous only on a \\emph{comeager} set, see the work of Bernshteyn \\cite{Bernshteyn2021LLL}.\nThe class $\\mathsf{MEASURE}$ generalizes the setting of fiid.\nThe goal here is to find a Borel solution on a conull set for a given Borel graph that is endowed with a Borel probability measure.\nThis direction has been recently investigated by Bernshteyn \\cite{Bernshteyn2021LLL} for general graphs, where he proved a measurable version of Lovasz's Local Lemma.\nSee \\cref{fig:big_picture_grids} for all the relations that are known to us.\n\\begin{figure}\n \\centering\n \\includegraphics[width=\\textwidth]{fig\/big_picture_grids.png}\n \\caption{Big picture of lcl for $\\mathbb{Z}^d$}\n \n\\end{figure}\n\n\n\n\n\n\n\n\\subsection{The Results}\n\nIn this subsection we formulate our main results.\nThe formulation uses mostly the language of ffiid because the tail behavior of the coding radius allows to compare problems on a quantitative and, as it turns out, rich scale.\nOne should always re-interpret the result with the help of Figure~\\ref{fig:big_picture_grids} to different contexts.\nIronically enough, the techniques that are most valuable when investigating ffiid solutions come either from local algorithms or Borel combinatorics, where such a rich and quantitative distinction is not possible.\nRecall that in the distributed algorithms there is a trichotomy of LCL problems, i.e., trivial, local or global, and in Borel combinatorics we distinguish mainly between trivial, continuous, Borel, and global problems.\n\n\nWe start with results that use ideas from distributed algorithms.\nFirst, we prove that LCL problems that are not global in distributed algorithms are precisely problems that admit ffiid solution with a tower tail decay.\n\n\\begin{theorem*}\nLet $d\\in \\mathbb{N}$ and $\\Pi$ be an LCL problem on $\\mathbb{Z}^d$.\nThen $\\Pi$ is not global in the local world, i.e., there is a deterministic algorithm in the local model that runs in $O(\\log^*(n))$ rounds, if and only if there is an ffiid solution $F$ of $\\Pi$ such that\n$$\\P(R_\\mathcal{A}>r)\\le \\frac{1}{tower(Cr)}$$\nfor some $C>0$ and large enough $r\\in \\mathbb{N}$.\n\\end{theorem*}\n\nThe translation between ffiid and distributed algorithms allows to deduce the following result that answers problem (ii) of Holroyd at el.~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\n\n\\begin{theorem*}\nLet $d\\in \\mathbb{N}$ and $\\Pi$ be an LCL problem on $\\mathbb{Z}^d$.\nIf $\\Pi$ admits an ffiid solution $F$ that has faster tail decay than any tower function, then $\\Pi$ is trivial, i.e., there is a color such that a constant coloring with this color is a correct solution.\n\\end{theorem*}\n\nAs a last application of ideas from the distributed algorithms, we show the following speed-up theorem for ffiid.\n\n\\begin{theorem*}\nLet $d,t\\in \\mathbb{N}$.\nThere is $C_{d,t}>0$ such that if an LCL problem $\\Pi$ admits an ffiid solution $F$ such that\n$$\\P(R_\\mathcal{A}>r)\\le C_{d,t}r^{d-1},$$\nthen there is a deterministic algorithm in the local model that runs in $O(\\log^*(n))$ rounds and solves $\\Pi$.\n\\end{theorem*}\n\nTo understand better LCL problems that are global in the distributed algorithms or admit no continuous solution in the context of Borel combinatorics we introduce the following numerical invariant.\nThe inspiration comes from the paper of Holroyd et al.~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\n\n\\begin{definition*}\nLet $\\Pi$ be an LCL on $\\mathbb{Z}^d$.\nWe define ${\\bf M}(\\Pi)\\in [0,+\\infty]$ to be the supremum over all $\\alpha\\in [0,+\\infty)$ such that there is an ffiid solution $F$ of $\\Pi$ such that $R_F$ has finite the $\\alpha$-moment.\n\\end{definition*}\n\nCombination of previous results immediately implies the following strengthening of the upper bound on the necessary infinite moment of ffiid solution of the proper vertex $3$-coloring problem, $\\Pi^d_{\\chi,3}$, from Holroyd et al.~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\nMoreover, this answers their problem (iii).\n\n\\begin{theorem*}\nLet $d\\in \\mathbb{N}$ and $\\Pi$ be an LCL problem on $\\mathbb{Z}^d$.\nIf ${\\bf M}(\\Pi)>d-1$, then $\\Pi$ admits an ffiid solution with tower tail decay.\n\\end{theorem*}\n\nAs a first application of ideas from Borel combinatorics, i.e., the toast structure, we compute the optimal tail decay for $\\Pi^d_{\\chi,3}$, thus answering a problem (i) of Holroyd et al.~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\n\n\\begin{theorem*}\nLet $d>1$ and $\\Pi^d_{\\chi,3}$ be the proper vertex $3$-coloring problem.\nThen ${\\bf M}(\\Pi^d_{\\chi,3})=1$.\n\\end{theorem*}\n\nA high level explanation of a toast construction in the context of ffiid says that if $r\\in \\mathbb{N}$ is given, then the fraction of vertices not covered by first $r$-layers of the toast is at most $\\frac{1}{r}$.\nSimilar construction that uses only $1$-dimensional skeleton of some rectangular tiling of $\\mathbb{Z}^d$ improve this upper bound to $\\frac{1}{r^{d-1}}$.\nWe show that this type of structure is enough to find a perfect matching in a bigger graph then $\\mathbb{Z}^d$, a graph where we moreover add diagonal edges in every $2$-dimensional hyperplane.\nThis yields an existence of intermediate problems from the point of view of ${\\bf M}$.\n\n\\begin{theorem*}\nLet $2\\le \\ell\\le d\\in \\mathbb{N}$.\nThen there is an LCL problem $\\Pi$ on $\\mathbb{Z}^d$ that satisfies ${\\bf M}(\\Pi)=\\ell-1$.\n\\end{theorem*}\n\nAs a last result we summarize what is known about specific problems.\nSince we work mainly with oriented grids we state everything in this context, however, some of the results can be deduced for non-oriented grids as well, see Section~\\ref{sec:SpecificProblems}.\nIn particular, we answer Question~6.5 of Gao et al.~\\cite{GJKS} about Borel edge chroamtic number of $\\operatorname{Free}(2^{\\mathbb{Z}^d})$.\n\n\\begin{theorem*}\nLet $d\\ge 2$ and $\\Pi$ be any of the following LCL problems:\n\\begin{enumerate}\n \\item proper vertex $3$-coloring, \\cite{GJKS}, \\cite{spencer_personal},\n \\item proper edge $2d$-coloring,\n \\item rectangular tiling with rectangles from some set $\\mathcal{R}$ that satisfies \\cref{eq:GCD}, \\cite{Unger},\n \\item perfect matching, \\cite{GJKS}.\n\\end{enumerate}\nThen $\\Pi$ is in the class $\\TOAST$.\nConsequently, $\\Pi$ is in the class $\\mathsf{BOREL}$ and satisfies $1\\le {\\bf M}(\\Pi)$.\n\\end{theorem*}\n\n\n\nThe paper is organized as follows...\n\n\n\\end{comment}\n\n\n\n\n\n\n\\section{Preliminaries}\n\\label{sec:preliminaries}\n\nLet $d\\in \\mathbb{N}$ be fixed.\nThe class of graphs that we consider in this work contains mostly the infinite grid $\\mathbb{Z}^d$, that is $\\mathbb{Z}^d$ with the Cayley graph structure given by the standard set of generators \n$$S=\\{(1,0,\\dots, 0),(0,1,\\dots, 0), \\dots (0,\\dots,0,1)\\}. $$\nWe also used the finite variant of the graph: the torus of side length $\\Theta(n^{1\/d})$ with edges oriented and labeled by $S$.\nUsually, when we talk about a graph $G$ we mean one of the graphs above and it follows from the context whether we mean finite or infinite graph. \n\nGiven a graph $G$, we denote as $V(G)$ the set of vertices, as $E(G)$ the set of edges and as $d_G$ the \\emph{graph distance} on $G$, i.e., $d_G(v,w)$ is the minimal number of edges that connects $v$ and $w$.\nWe write $\\mathcal{B}_G(v,r)$ for the \\emph{ball of radius $r\\in \\mathbb{N}$} around $v\\in V(G)$. When $G$ is the infinite grid $\\mathbb{Z}^d$, then we leave the subscripts blank.\nA \\emph{rooted graph} is a pair $(G,v)$ where $G$ is a graph and $v\\in V(G)$.\nWe denote as $[G,v]$ the \\emph{isomorphism type} of the rooted graph $(G,v)$, where $(G,v)$ and $(G',v')$ are isomorphic if there is a bijection $\\varphi:V(G)\\to V(G')$ that sends $v$ to $v'$ and preserves edge orientation and labeling.\n\nLet $D$ be a set, $G$ a graph and $c:V(G)\\to D$ be a labeling of vertices.\nThen we write $G_c=(G,c)$ for the pair graph + labeling.\nSimilarly we denote as $(G_c,v)$ the rooted graph $(G,v)$ with the labeling $c$, $[G_c,v]$ its isomorphism class and $\\mathcal{B}_{G_c}(v,r)$ the labeled ball of radius $r\\in \\mathbb{N}$ around $v$, where $v\\in V(G)$.\nOften we do not mention the labels and tacitly assume that they are understood from the context.\nAlso, the label set might depend on the size of of the vertex set, e.g., unique labels from $[|V(G)|]=\\{1,\\dots,|V(G)|\\}$ etc.\n\n\\begin{comment}\n\nWe always consider the group $\\mathbb{Z}^d$ with the standard set of generators $S$, i.e.,\n$$S=\\{(1,0,\\dots, 0),(0,1,\\dots, 0), \\dots (0,\\dots,0,1)\\}.$$\nThe \\emph{Cayley graph} $\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right):=\\operatorname{Cay}\\left(\\mathbb{Z}^d,S\\right)$ of $\\mathbb{Z}^d$ is then defined on the vertex set $\\mathbb{Z}^d$ with oriented and $S$-labeled edges \n$$\\left\\{(g,s+g):s\\in S, g\\in \\mathbb{Z}^d\\right\\}=E\\left(\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)\\right).$$\nWe denote the identity in $\\mathbb{Z}^d$ as ${\\bf 0}_d$. \nWhen we talk about $\\mathbb{Z}^d$ next, we mean formally the appropriate Cayley graph with labelled edges, that is, we think of the graph together with its consistent orientation. \n\n\nA \\emph{graph} $G$ is a pair $(V(G),E(G))$, where $V(G)$ is a vertex set and $E(G)$ consists of oriented and $S$-labeled edges.\nGiven a graph $G$, we denote as $d_G$ the \\emph{graph distance} on $G$, i.e., $d_G(v,w)$ is the minimal number of edges that connects $v$ and $w$.\nWe write $\\mathcal{B}_G(v,r)$ for the \\emph{ball of radius $r\\in \\mathbb{N}$} around $v\\in V(G)$. Usually, $G$ is just the grid $\\mathbb{Z}^d$, then we leave the subscript blank. \nA \\emph{rooted graph} is a pair $(G,v)$ where $G$ is a graph and $v\\in V(G)$.\nWe denote as $[G,v]$ the isomorphism type of the rooted graph $(G,v)$, where $(G,v)$ and $(G',v')$ are isomorphic if there is a bijection $\\varphi:V(G)\\to V(G')$ that sends $v$ to $v'$ and preserves edge orientation and labeling.\n\n\\todo{vsechny strany toru maji delku n na 1\/d; zjevne vsechny problemy maji lokalitu nejvyse n na 1\/d}\n\n\\end{comment}\n\n\n\\begin{comment}\nWe write $\\operatorname{\\bf FG}^{S}$ for the set of all graphs that look in their $2$-neighborhoods as $\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)$, i.e., a graph $G=(V(G),E(G))$ is in $\\operatorname{\\bf FG}^{S}$ if $[\\mathcal{B}_G(v,2)]=\\left[\\mathcal{B}_{\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)}({\\bf 0}_d,2)\\right]$ for every $v\\in V(G)$.\nNote that $\\mathbb{Z}^d\\in \\operatorname{\\bf FG}^{S}$ and there are finite graphs in $\\operatorname{\\bf FG}^{S}$, as an example take any torus that is big enough.\nWrite $\\operatorname{\\bf FG}_{\\bullet}^{S}$ for the rooted graphs of the form\n$$(\\mathcal{B}_G(v,r),v)$$\nwhere $G\\in \\operatorname{\\bf FG}^{S}$, $v\\in V(G)$ and $r\\in \\mathbb{N}$.\nWe say that a rooted graph $(G,v)\\in \\operatorname{\\bf FG}_{\\bullet}^{S}$ has \\emph{diameter at most $m\\in \\mathbb{N}$} if $(\\mathcal{B}_G(v,m),v)=(G,v)$.\nNote that, in particular, we have $\\mathcal{B}_{\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)}({\\bf 0}_d,r)\\in \\operatorname{\\bf FG}_{\\bullet}^{S}$ for every $r\\in \\mathbb{N}$.\n\nLet $D$ be a set, $G$ a graph and $c:V(G)\\to D$ be a labeling of vertices.\nThen we write $G_c=(G,c)$ for the pair graph + labeling.\nSimilarly we denote as $(G_c,v)$ the rooted graph $(G,v)$ with the labeling $c$, $[G_c,v]$ its isomorphism class and $\\mathcal{B}_{G_c}(v,r)$ the labeled ball of radius $r\\in \\mathbb{N}$ around $v$, where $v\\in V(G)$.\nOften the label set $D$ depends on the size of of the vertex set, e.g., unique labels from $[|V(G)|]=\\{1,\\dots,|V(G)|\\}$ etc.\n\\end{comment}\n\n\n\n\\subsection{LCLs on $\\mathbb{Z}^d$}\n\nOur goal is to understand locally checkable problems (LCL) on $\\mathbb{Z}^d$. \n\n\\begin{definition}[Locally checkable problems \\cite{naorstockmeyer}]\nA \\emph{locally checkable problems (LCL)} on $\\mathbb{Z}^d$ is a quadruple $\\Pi=(\\Sigma_{in}, \\Sigma_{out},t,\\mathcal{P})$ where $\\Sigma_{in}$ and $\\Sigma_{out}$ are finite sets, $t\\in \\mathbb{N}$ and $\\mathcal{P}$ is an assignment\n$$[H_c,w]\\mapsto \\mathcal{P}([H_c,w])\\in \\{0,1\\},$$\nwhere $(H,w)$ is a rooted grid of radius at most $t$ together with an input and output coloring $c: V(H)\\to \\Sigma_{in}\\times\\Sigma_{out}$. \n \nGiven a graph $G$ together with a coloring $c:V(G)\\to \\Sigma_{in}\\times\\Sigma_{out}$, we say that $c$ is a \\emph{$\\Pi$-coloring} if\n$$\\mathcal{P}([\\mathcal{B}_{G_c}(v,t),v])=1$$\nfor every $v\\in V(G)$.\n\\end{definition}\n\n\n\n\\begin{comment}\n\nClasses of functions\n\nBefore we proceed to definition of local algorithms and their complexity we recall here the relevant notions concerning functions from $\\mathbb{N}$ to $\\mathbb{N}$. \nImportantly, we often use asymptotic notation as usual in Computer Science, e.g., we define the class of functions $O(f(n))$ as the union over all functions $g:\\mathbb{N} \\rightarrow \\mathbb{N}$ such that there exists $C > 0$ such that for all $n$ we have $g(n) \\le C f(n)$. Other symbols $\\Omega, \\Theta, o, \\omega$ are defined analogously and their usage is analogous (see e.g. \\cite{?}). \n\\end{comment}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{LOCAL model}\n\\label{subsec:preliminaries_local}\n\nHere we formally introduce the $\\mathsf{LOCAL}$ model for both deterministic and randomized algorithms and state the classification of LCLs on $\\mathbb{Z}^d$ from the $\\mathsf{LOCAL}$ point of view.\nInformally, a local algorithm $\\mathcal{A}$ of radius $r\\in \\mathbb{N}$ is a rule that takes a labeled graph $G_c$ and then re-labels each vertex $v\\in V(G)$ depending only on the isomorphism type of its neighborhood $[\\mathcal{B}_{G_c}(v,r),v]$ with a new label from some finite set $b$.\n\n\\paragraph{Local algorithms}\n\nHere, we formally define the complexity classes $\\mathsf{LOCAL}(\\mathcal{F})$ of deterministic local algorithms, where $\\mathcal{F}$ is a class of functions from $\\mathbb{N}$ to $\\mathbb{N}$. \nThe somewhat unintuitive part of the definition is that the algorithm takes as input \\emph{unique identifiers}. Usually, they are taken from some range of values that is polynomial in $n$, the number of vertices of the input graph. However, in most scenarios this is equivalent to the assumption that they are from range $\\{1,\\dots n\\}$, i.e., every identifier is used exactly once. We also use this definition here, for simplicity.\nIn all forthcoming definition we have some fixed LCL $\\Pi$ in mind.\n\n\\begin{definition}[A deterministic local algorithm for $\\mathbb{Z}^d$]\n\\label{def:local_algorithm\nLet $\\Sigma_{out}$ be a finite set.\nA \\emph{local (deterministic) algorithm $\\mathcal{A}$} is a sequence of pairs $\\{\\mathcal{A}_n,\\ell_n\\}_{n\\in\\mathbb{N}}$, where $\\ell_n\\in \\mathbb{N}$ and $\\mathcal{A}_n$ is an assignment, such that\n$$[\\mathcal{B}_{G_\\ell}(v,\\ell_n),v]\\mapsto \\mathcal{A}_n([\\mathcal{B}_{G_\\ell}(v,\\ell_n),v])\\in \\Sigma_{out}$$\nfor every $\\Sigma_{in}$-labeled graph $G$ (torus) on $n$ vertices, $v\\in V(G)$ and a bijection $\\ell:V(G)\\to [n]$.\nIn another words, $\\mathcal{A}_n$ is a function that is defined on isomorphism classes of rooted graphs of diameter at most $\\ell_n$ with injectively labeled vertices from the set $\\{1,\\dots, n\\}$ and with range $b$. \n\nGiven such a $G$ and a bijection $\\ell:V(G)\\to |V(G)|$, we define\n$$\\mathcal{A}(G_\\ell)(v):=\\mathcal{A}_{|V(G)|}([\\mathcal{B}_{G_\\ell}(v,\\ell_n),v])$$\nfor every $v\\in V(G)$.\n\\end{definition}\n\nNext we define the deterministic complexity of a given local coloring problem $\\Pi$.\n\n\\begin{definition}[Complexity of deterministic local algorithm for $\\mathbb{Z}^d$]\nWe say that a local algorithm $\\mathcal{A}=\\{\\mathcal{A}_n,\\ell_n\\}_{n\\in \\mathbb{N}}$ \\emph{solves} an LCL $\\Pi=(\\Sigma_{in},\\Sigma_{out},t,\\mathcal{P})$ if $\\mathcal{A}(G_l)$ is a $\\Pi$-coloring of $G$ for every finite $\\Sigma_{in}$-labeled graph $G$ and any bijection $\\ell:V(G)\\to |V(G)|$.\n\nWe say that $\\Pi$ is in the class $\\mathsf{LOCAL}(\\mathcal{F})$, where $\\mathcal{F}$ is a class of functions, if $n\\mapsto \\ell_n$ is in $\\mathcal{F}$.\n\\end{definition}\n\nA randomized local algorithm and its complexity is defined analogously but instead of labeling with unique identifiers each node has an iid random string and the algorithm succeeds with probability $1 - 1\/n$.\n\nCrucial for our investigation are so-called \\emph{oblivious local algorithms}.\nRoughly speaking an algorithm is \\emph{oblivious} if it does not have knowledge of $n$ (see \\cite{Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring} for details). \nSince we do not consider deterministic oblivious algorithms, by an oblivious local algorithm we mean an oblivious randomized local algorithm.\n\n\\begin{definition}[Oblivious local algorithm]\n\\label{def:uniform_local}\nA \\emph{oblivious local algorithm} $\\mathcal{A}$ is a pair $(\\mathcal{C},f)$, where $\\mathcal{C}$ is a (measurable) collection of rooted graphs labeled with $\\Sigma_{in}$ and real numbers from $[0,1]$, and $f: \\mathcal{C} \\rightarrow \\Sigma_{out}$ is a measurable assignment.\n\nWe define $\\mathcal{A}(G)$ to be the labeling computed by $\\mathcal{A}$ on $G$, where for a given node $u$ we have $\\mathcal{A}(G)(u)=g([B_G(u, t), u])$ where $t$ is the smallest integer such that $[B_G(u, t),u] \\in \\mathcal{C}$.\n\\end{definition}\n\nNote that since there is no $n$ in the definition of oblivious local algorithm, the algorithm is correct if it always finishes with probability $1$, while randomized algorithms can have small probability of error dependent on $n$. Hence, one should think about such algorithms as an analogue of Las Vegas algorithms, while the classical randomized local algorithms correspond to Monte Carlo algorithms. \nAn advantage of not knowing $n$ is that it is possible to run $\\mathcal{A}$ on, e.g., the infinite grid $\\mathbb{Z}^d$ instead of just a torus. Usually, these two models give the same results but for toast algorithms that we consider later, there is a difference: While we define the algorithms for an infinite grid, for example the $3$-coloring algorithm makes sense also for a finite torus, while the perfect matching algorithm does not make sense in the finite world for the simple reason that the input graph can have an odd number of vertices. \n\n\n\n\n\n\n\n\n\\begin{comment}\n\\paragraph{Randomized local algorithms}\\todo{merge with previous paragraph}\n\nThe main difference between deterministic and randomized model is that the input of the local rule in the randomized case are graphs with vertices labeled uniformly independently at random by elements of the interval $[0,1]$ and we require that the output coloring is correct everywhere with high probability.\nThe complexity of local algorithm is then defined as the locality it needs.\n\n\\begin{definition}[Randomized local algorithm for $\\mathbb{Z}^d$]\nLet $b$ be a set.\nWe say that $\\mathbb{A}=\\{\\mathbb{A}_n,\\ell_n\\}_{n\\in\\mathbb{N}}$ is a randomized local algorithm for the class $\\operatorname{\\bf FG}^{S}$ if $\\ell_n\\in \\mathbb{N}$ and \n$$[\\mathcal{B}_{G_l}(v,\\ell_n),v]\\mapsto \\mathbb{A}_n([\\mathcal{B}_{G_l}(v,\\ell_n),v])\\in b$$\nis a measurable assignment, where $G\\in \\operatorname{\\bf FG}^{S}$, $|V(G)|=n$, $v\\in V(G)$ and $\\ell:V(G)\\to [0,1]$.\nIn another words, $\\mathbb{A}_n$ is a measurable function that is defined on isomorphism classes of rooted graphs of diameter at most $\\ell_n$ with vertices labeled by $[0,1]$ and with range $b$. \n\nGiven $G\\in \\operatorname{\\bf FG}^{S}$ such that $n=|V(G)|$ and a map $\\ell:V(G)\\to [0,1]$, we define\n$$\\mathbb{A}(G_l)(v):=\\mathbb{A}_n([\\mathcal{B}_{G_l}(v,\\ell_n),v])$$\nfor every $v\\in V(G)$.\n\\end{definition}\n\nNext we define when a randomized algorithm solves an LCL $\\Pi$ and the randomized complexity classes.\n\n\\begin{definition}[Complexity of randomized local algorithm for $\\mathbb{Z}^d$]\nLet $\\Pi=(b,t,\\mathcal{P})$ be an LCL.\nWe say that a randomized algorithm $\\mathbb{A}=\\{\\mathbb{A}_n,\\ell_n\\}_{n\\in \\mathbb{N}}$ solves $\\Pi$ if \n$$\\lambda^{\\oplus n}\\left(\\left\\{\\ell\\in [0,1]^{V(G)}:\\mathbb{A}(G_l) \\ \\operatorname{is \\ a} \\ \\Pi\\operatorname{-coloring} \\right\\}\\right)>1-\\frac{1}{n}$$\nholds for every $n\\in \\mathbb{N}$ and $G\\in \\operatorname{\\bf FG}^{S}$ such that $|V(G)|=n$ and\n$$[\\mathcal{B}_{G}(v,\\ell_n+t),v]=\\left[\\mathcal{B}({\\bf 0}_d,\\ell_n+t),{\\bf 0}_d\\right]$$\nfor every $v\\in V(G)$.\n\nWe say that $\\Pi$ is in the class $\\mathsf{R-LOCAL}(\\mathcal{F})$, where $\\mathcal{F}$ is a class of functions, if $n\\mapsto \\ell_n$ is in $\\mathcal{F}$.\n\\end{definition}\n\\end{comment}\n\n\\begin{comment}\n\\label{sec:clocal_definition}\n\nThe last class of local algorithm that we introduce serves on one hand as a main technical tool for deterministic speed up tricks and on the other hand allows a smooth translation of local algorithm into the measurable setting.\nIntuitively, an algorithm is in the class $\\mathsf{C-LOCAL}$ if it only needs a fixed finite locality (not depending on the size of a graph) given some coloring of the vertices with fixed number of colors (not depending on the size of a graph).\n\nLet $s\\in \\mathbb{N}$ and $G$ be a graph.\nDenote as $G^s$ the $s$-power graph of $G$, i.e., $V(G^s)=V(G)$ and $(x,y)\\in E(G^s)$ if and only if $0 r),$$\nwhat we mean is the following. \nFirst, for a given rooted neighborhood $(G,v)$ with a fixed $\\Sigma_{in}$-labeling we ask what is the probability that, after labeling $G$ with real numbers uniformly at random, that $(G,v)$ is in $\\mathcal{C}$. \nSecond, we take the supremum over all such values over all $\\Sigma_{in}$-labels.\n\n\\begin{definition}[$\\mathsf{TAIL}(f)$]\n\\label{def:tail}\nWe say that an LCL $\\Pi$ is in the class $\\mathsf{TAIL}(f)$ for some function $f:\\mathbb{N}\\to \\mathbb{N}$ if there is an oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi$ such that\n$$\\P(R_\\mathcal{A} > r)\\le \\frac{1}{f(r)}$$\nfor every large enough $r\\in \\mathbb{N}$. \n\nWe say that $\\Pi \\in \\mathsf{TAIL}$ if there exists $f \\in \\omega(1)$ such that $\\Pi \\in \\mathsf{TAIL}(f)$. \nWe say that $\\Pi$ is solved by a \\emph{block factor} if there is an oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi$ such that for some $r_0$ we have $\\P(R_\\mathcal{A} > r_0) = 0$.\n\\end{definition}\n\nHere things are made simple because we think of the algorithm as being run on an infinite graph so nothing depends on $n$ in this complexity measure (while it is a crucial parameter in the classical local complexity definition). \n\nWe also use the asymptotic notation when talking about $\\mathsf{TAIL}$ complexities. For example, an important class $\\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$ is defined as the class of LCLs $\\Pi$ such that there is a constant $C\\in \\mathbb{N}$ and oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi$ that satisfies\n$$\\P(R_\\mathcal{A}> r)\\le \\frac{1}{\\mathrm{tower}(Cr)}$$\nwhere $\\mathrm{tower}(x)=\\exp\\circ\\dots\\circ \\exp(1)$ is $\\lfloor x\\rfloor$-many times iterated exponential.\n\nA coarser way to measure the complexity of a problem $\\Pi$ is to ask what moments of $R_\\mathcal{A}$ could be finite, when $\\mathcal{A}$ is an oblivious local algorithm.\nFor example, the first moment finite corresponds to the fact that average time until an algorithm finishes is defined (for finite graphs this is independent on the size of the graph). \nIn particular, we investigate the question what is the least infinite moment for any oblivious local algorithm to solve a given problem $\\Pi$, as this language was used in \\cite{HolroydSchrammWilson2017FinitaryColoring}. \nRecall that $R_\\mathcal{A}$ has a \\emph{finite $\\alpha$-moment}, where $\\alpha\\in [0,+\\infty)$, if\n$$\\mathbb{E}\\left(\\left(R_\\mathcal{A}\\right)^\\alpha\\right)=\\sum_{r \\in \\mathbb{N}} r^\\alpha \\cdot \\P(R_\\mathcal{A} = r) <+\\infty.$$\n\n\\begin{definition}\nLet $\\Pi$ be an LCL.\nWe define ${\\bf M}(\\Pi)\\in [0,+\\infty]$ to be the supremum over all $\\alpha\\in [0,+\\infty)$ such that there is an oblivious $\\mathcal{A}$ that solves $\\Pi$ and has the $\\alpha$-moment of $R_\\mathcal{A}$ finite.\n\\end{definition}\n\nThe $\\alpha$ moment being finite roughly corresponds to $\\Omega(r^{\\alpha})$ $\\mathsf{TAIL}$ complexity. In fact we now recall the following observation.\n\n\\begin{fact}\n\\label{fact:moment_implies_tail}\nSuppose that the coding radius of an oblivious local algorithm $\\mathcal{A}$ has the $\\alpha$-moment finite.\nThen for every $C>0$ we have\n$$\\P(R_\\mathcal{A} > r) \\le C\\frac{1}{r^\\alpha}$$\nfor large enough $r\\in \\mathbb{N}$.\n\\end{fact}\n\\begin{proof}\nWe prove the contrapositive. Suppose the tail of the coding radius $R$ satisfies $\\P(R > r_i) > C \\cdot \\frac{1}{r^\\alpha_i}$ for some $C > 0$ and an infinite sequence $(r_i)_{i\\in \\mathbb{N}}$. \nAs the sequence $\\P(R> r_i)$ goes to zero, we may pass to a subsequence such that for each $i\\in \\mathbb{N}$ we have $\\P(R > r_{i+1}) \\le \\frac{1}{2}\\P(R > r_i)$. We call the new subsequence still $(r_i)_{i\\in \\mathbb{N}}$. \nWe have \n\\begin{align*}\n\\mathbb{E}(R^\\alpha) &\\ge \\sum_{i \\in \\mathbb{N}} r_i^\\alpha \\cdot (\\P(R> r_i) - \\P(R > r_{i+1}))\\\\\n&\\ge \\sum_{i \\in \\mathbb{N}} r_i^\\alpha \\cdot \n\\frac{1}{2}\\P(R > r_i)\n\\ge \\sum_{i \\in \\mathbb{N}} r_i^\\alpha \\cdot \n\\frac{C}{2} \\cdot \\frac{1}{r_i^\\alpha}\n= +\\infty\n\\end{align*}\nwhich finishes the proof. \n\\end{proof}\n\nFinally, an important fact is that a small $\\mathsf{TAIL}$ complexity of an oblivious randomized local algorithm implies that it has small randomized complexity (on finite graphs). This is because if we run an oblivious algorithm with tail decay $1\/f(r)$ on an $n$-node graph and mark its run as a failure if the coding radius exceeds $f^{-1}(n^2)$, the failure at some node happens only with probability $1\/n$. This was formally proven in \\cite{brandt_grebik_grunau_rozhon2021classification_of_lcls_trees_and_grids} and here we only state the result. \n\n\\begin{lemma}\n\\label{lem:tail->local}\nLet $\\mathcal{A}$ be an oblivious local randomized algorithm with $\\mathsf{TAIL}$ complexity $f(r)$. Then the randomized complexity of $\\mathcal{A}$ is $O(f^{-1}(n^2))$. \n\\end{lemma}\n\nIn fact, one can prove that for $o(\\log n)$ local complexities the two notions are the same. We refer to the reader to the discussion in \\cite{brandt_grebik_grunau_rozhon2021classification_of_lcls_trees_and_grids} for more details. \n\n\\begin{theorem}[\\cite{brandt_grebik_grunau_rozhon2021classification_of_lcls_trees_and_grids}]\n\\label{thm:tail=uniform}\nLet $\\mathcal{A}$ be an oblivious local randomized algorithm with randomized complexity $f(n) = o(\\log n)$. Then this algorithm in fact never returns an incorrect solution and its tail complexity is $\\Omega(f^{-1}(r))$. \n\\end{theorem}\n\n\n\n\n\n\n\n\n\n\\subsection{Subshifts of finite type}\n\\label{subsec:subshifts}\n\nHere we connect our definition of LCLs and oblivious local algorithms with the well-studied notions of subshifts of finite type and finitary factors of iid labelings.\n\nTo define subshifts of finite type we restrict ourselves to LCLs without input.\nThat is $\\Sigma_{in}$ is empty or contains only one element.\nMoreover, for convenience we set $b:=\\Sigma_{out}$.\nThen an LCL $\\Pi$ is of the form $\\Pi=(b,t,\\mathcal{P})$.\nWe define $X_\\Pi\\subseteq b^{\\mathbb{Z}^d}$ to be the space of all $\\Pi$-colorings.\nThat is we consider all possible labelings of $\\mathbb{Z}^d$ by elements from $b$ and then include in $X_{\\Pi}$ only those that satisfies the finite constraints $\\mathcal{P}$.\nViewing $\\mathbb{Z}^d$ as a group, there is a natural shift action of $\\mathbb{Z}^d$ on (the Cayley graph) $\\mathbb{Z}^d$.\nNamely, $g\\in \\mathbb{Z}^d$ shifts $h$ to $g+h$.\nThis action naturally ascend to an action $\\cdot $ on $b^{\\mathbb{Z}^d}$ as\n$$(g\\cdot x)(h)=x(g+h)$$\nwhere $x\\in b^{\\mathbb{Z}^d}$.\nIt is not hard to see (since $\\mathcal{P}$ is defined on isomorphism types and not rooted in any particular vertex) that $X_\\Pi$ is invariant under the shift map.\nThen we say that $X_\\Pi$ together with the shift structure is a \\emph{subshift of finite type}.\n\nA \\emph{finitary factor of iid process (ffiid) $F$} is a function that takes as an input $\\mathbb{Z}^d$ labeled with real numbers from $[0,1]$, independently according to Lebesgue measure, and outputs a different labeling of $\\mathbb{Z}^d$, say with elements of some set $b$. \nThis labeling satisfies the following: (a) $F$ is equivariant, that is for every $g\\in \\mathbb{Z}^d$ shifting the real labels with $g$ and applying $F$ is the same as applying $F$ and then shifting with $g$, (b) the value of $F(x)({\\bf 0})$ depends almost surely on the real labels on some random but finite neighborhood of ${\\bf 0}$.\nSimilarly as for oblivious local algorithms one defines the random variable $R_F$ that encodes minimal such neighborhood.\nWe say that ffiid $F$ solves $\\Pi$ if $F(x)\\in X_{\\Pi}$ with probability $1$.\n\nThe difference between ffiid and oblivious local algorithms is just formal.\nNamely, given an oblivious local algorithm $\\mathcal{A}$, we can turn it into an ffiid $F_\\mathcal{A}$ simply by running it simultaneously in every vertex of $\\mathbb{Z}^d$.\nWith probability $1$ this produces a labeling of $\\mathbb{Z}^d$ and it is easy to see that both conditions from the definition of ffiid are satisfied.\nOn the other hand, if $F$ is ffiid, then we collect the finite neighborhoods that determine the value at ${\\bf 0}$ together with the decision about the label.\nThis defines a pair $\\mathcal{A}_F=(\\mathcal{C},f)$ that satisfies the definition of an oblivious local algorithm.\nIt follows that the coding radius $R_\\mathcal{A}$, resp $R_F$, is the same after the translations.\nWe chose to work with the language of distributed algorithms as we borrow the techniques mainly from that area. \n\n\\begin{remark}\nIf in (b), in the definition of ffiid, we only demand that the value at ${\\bf 0}$ depends measurably on the labeling, i.e., $F$ is allowed to see the labels at the whole grid, then we say that $F$ is a \\emph{factor of iid labels (fiid)}.\n\nGiven an LCL $\\Pi$, we say that $\\Pi$ is in the class $\\mathsf{fiid}$, if there is a fiid $F$ that produces element of $X_\\Pi$ almost surely.\nIt is easy to see that $\\mathsf{TAIL}\\subseteq \\mathsf{fiid}$.\nWe remark that in the setting without inputs we have $\\mathsf{TAIL}=\\mathsf{fiid}$ on $\\mathbb{Z}$ and it is unknown, however, still plausible, whether this holds for $d>1$.\nFor general LCLs with inputs (where the definition of $\\mathsf{fiid}$ is extended appropriately) we have $\\mathsf{TAIL}\\not=\\mathsf{fiid}$ even for $d=1$, see \\cite{grebik_rozhon2021LCL_on_paths}.\n\\end{remark}\n\n\\begin{comment}\nGiven an LCL $\\Pi=(b,t,\\mathcal{P})$, we define $X_\\Pi\\subseteq b^{\\mathbb{Z}^d}$ to be the space of all $\\Pi$-colorings.\nWe say that an LCL $\\Pi$ is \\emph{satisfiable} if $\\emptyset\\not=X_\\Pi$, i.e., if there is a $\\Pi$-coloring of $\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)$.\nWe say that an LCL $\\Pi$ is \\emph{trivial} if there is $\\alpha\\in b$ such that the constant function $\\alpha$ is a $\\Pi$-coloring.\n\nThere is a canonical (right-)shift action $\\cdot$ of $\\mathbb{Z}^d$ on $b^{\\mathbb{Z}^d}$ (and in general on any $X^{\\mathbb{Z}^d}$, where $X$ is a set) that is defined as\n$$(g\\cdot x)(h)=x(h + g)$$\nfor every $g,h\\in \\mathbb{Z}^d$ and $x\\in b^{\\mathbb{Z}^d}$.\nIt is easy to see that $X_\\Pi$ is invariant under this action, i.e., $x\\in X_\\Pi$ implies $g\\cdot x\\in X_\\Pi$ for every $g\\in \\mathbb{Z}^d$.\nMoreover, since $b$ is a finite set, the product topology turns $b^\\Gamma$ into a compact metric space and it is easy to see $X_\\Pi$ is closed subspace.\nThe set $X_\\Pi$ together with the shift action is called a \\emph{sub-shift of finite type}.\n\nSuppose that $\\mathbb{Z}^d$ has an action $\\star_Y$ and $\\star_Z$ on spaces $Y$ and $Z$, respectively. \nWe say that a function $f:Y\\to Z$ is \\emph{equivariant} if $f(g\\star_Y x)=g\\star_Z f(x)$ for every $g\\in \\mathbb{Z}^d$ and $x\\in Y$.\n\nGiven an LCL $\\Pi$, our aim is to study complexity of finitary factors of iid that produce elements in $X_\\Pi$ almost surely.\nWe say that $\\Pi$ admits a \\emph{fiid solution} if there is a measurable equivariant map $F:[0,1]^{\\mathbb{Z}^d}\\to X_\\Pi$ that is defined on a $\\mu$-conull set, where $\\mu$ is the product measure on $[0,1]^{\\mathbb{Z}^d}$.\n\nWrite $\\varphi_F:[0,1]^{\\mathbb{Z}^d}\\to b$ for the composition of $F$ with the projection to ${\\bf 0}_d$-coordinate, i.e., \n$$x\\mapsto \\varphi_F(x)=F(x)({\\bf 0}_d).$$\nWe say that $F$ is a \\emph{finitary fiid (ffiid)} if for $\\mu$-almost every $x\\in [0,1]^{\\mathbb{Z}^d}$ the value of $\\varphi_F$ depends only on values of $x$ on some finite neighborhood of ${\\bf 0}_d$.\nThat is to say, there is a random variable $R_F$ that takes values in $\\mathbb{N}$, is defined $\\mu$-almost everywhere and satisfies for $\\mu$-almost every $x\\in [0,1]^{\\mathbb{Z}^d}$ the following: if $y\\in [0,1]^{\\mathbb{Z}^d}$ is such that\n$$x\\upharpoonright \\mathcal{B}_{\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)}\\left({\\bf 0}_d,R_\\varphi(x)\\right)=y\\upharpoonright \\mathcal{B}\\left({\\bf 0}_d,R_\\varphi(x)\\right),$$\nthen $\\varphi_F(x)=\\varphi_F(y)$.\n\\todo{nejak to zakoncit tim ze to nebudeme potrebovat}\n\\end{comment}\n\n\n\n\n\\begin{comment}\n\nWarm-up: Equivalence of $\\mathsf{LOCAL}(\\log^* n)$, $\\mathsf{TAIL}(\\mathrm{tower})$, $\\mathsf{CONTINUOUS}$\n\nRecall that $\\mathsf{LOCAL}(O(\\log^* n))$ denotes the class of satisfiable LCLs that satisfy (1) or (2) in Theorem~\\ref{thm:classification_of_local_problems_on_grids}.\nIn this section we show that every problem from $\\mathsf{LOCAL}(O(\\log^* n))$ is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$.\nIn the proof we use the fact that $\\mathsf{LOCAL}=\\mathsf{C-LOCAL}(O(1))$, see~Theorem~\\ref{thm:classification_of_local_problems_on_grids}.\n\nWe let $\\Delta_{s}:=|\\mathcal{B}({\\bf 0}_d,s)|$, where $s\\in \\mathbb{N}$.\nDenote as $\\mathcal{M}_s$ the subspace of $(\\Delta_s+1)^{\\mathbb{Z}^d}$ that consists of all proper vertex $(\\Delta_s+1)$-colorings of the $s$-power graph of $\\mathbb{Z}^d$.\nIt is easy to see that $\\mathcal{M}_s$ is nonempty, $\\mathbb{Z}^d$-invariant and compact.\nIn fact it is a subshift of finite type.\nThe following result holds for any graph of bounded degree, as was proved in \\cite{HolroydSchrammWilson2017FinitaryColoring}.\n\n\\begin{theorem}\\cite[Theorem~6]{HolroydSchrammWilson2017FinitaryColoring}\\label{th:ColorTower}\nLet $s\\in\\mathbb{N}$.\nThen there is an ffiid $F:[0,1]^{\\mathbb{Z}^d}\\to \\mathcal{M}_s$ and a constant $C_s>0$ such that\n$$\\mu(R_F>r)<\\frac{1}{tower(C_sr)}.$$\nIn particular, the LCL of coloring the $s$-power graph of $\\mathbb{Z}^d$ with $(\\Delta_s+1)$-many colors is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$.\n\\end{theorem}\n\n\\begin{proposition}\\label{pr:LOCALandCOLORING}\nLet $\\Pi=(b,t,\\mathcal{P})$ be a satisfiable LCL problem that is in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\nThen there is a continuous equivariant function $\\varphi:\\mathcal{M}_s\\to \\mathcal{C}_{\\Pi}$ for some $s\\in \\mathbb{N}$.\n\\end{proposition}\n\\begin{proof}\nSuppose that $\\Pi$ is in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\nBy Theorem~\\ref{thm:classification_of_local_problems_on_grids}, we have that $\\Pi$ is in the class $\\mathsf{C-LOCAL}(O(1))$.\nFix $r,\\ell\\in \\mathbb{N}$ and $\\mathcal{A}$ that witnesses this fact.\nLet $x\\in \\mathcal{M}_{r+t}$.\nThen define\n$$\\varphi(x)(h):=\\mathcal{A}([\\mathcal{B}_{\\mathbb{Z}^d_x}(h,\\ell),h]).$$\nIt is clear that $\\varphi$ is well defined, because $\\ellr)<\\frac{\\Delta_{\\ell}}{tower(C_s(r-\\ell))}\\le \\frac{1}{tower(Dr)}$$\nfor $r\\in \\mathbb{N}$ large enough and some constant $D\\in (0,+\\infty)$.\n\\end{proof}\n\n\n\n\n\ngives that\n\n\n\n\n\n\n\n\n\n$\\mathsf{TAIL}(\\tow(\\Omega(r)))$ $\\Rightarrow$ $\\mathsf{LOCAL}(O(\\log^* n))$\n\n\nFirst we show how to get randomized local algorithm from ffiid in a qualitative way.\nWe note that similar results holds for any bounded degree graph.\n\n\\begin{theorem}\\label{th:TailToRandom}\nLet $\\Pi$ be an LCL problem that is in the class $\\mathsf{TAIL}(f)$ for some $f:\\mathbb{N}\\to \\mathbb{N}$.\nSet $g(n)=\\min\\{m\\in \\mathbb{N}:f(m)\\ge n^2\\}$ and suppose that $g\\in o(n^{\\frac{1}{d}})$.\nThen $\\Pi$ is in the class $\\mathsf{R-LOCAL}(O(g(n)))$.\n\\end{theorem}\n\\begin{proof}\nFix an ffiid $F:[0,1]^{\\mathbb{Z}^d}\\to X_\\Pi$ such that\n$$\\mu(R_F>r)\\le \\frac{1}{f(r)}$$\nfor large enough $r\\in \\mathbb{N}$.\nDenote as $\\varphi_F$ the finitary rule and as $\\operatorname{dom}(\\varphi_F)$ its domain.\nThat is $\\alpha\\in \\operatorname{dom}(\\varphi_F)$ if $\\alpha\\in [0,1]^{V(\\mathcal{B}({\\bf 0}_d,s))}$ for some $s\\in \\mathbb{N}$ and satisfies\n$$F(x)=F(y)$$\nfor every extensions $\\alpha\\sqsubseteq x,y\\in [0,1]^{\\mathbb{Z}^d}$.\nIn that case we set $\\varphi_F(\\alpha)=F(x)$ for any such extension $\\alpha\\sqsubseteq x$.\n\nDefine $\\ell_n:=g(n)$.\nLet $(G,v)\\in \\operatorname{\\bf FG}_{\\bullet}^{S}$ be of a radius at most $\\ell_n$.\nConsider any map $\\ell:V(G)\\to [0,1]$.\nDefine\n$$\\mathbb{A}_n([G_l,v])=\\varphi_F(\\ell)$$\nin the case when $\\ell\\in \\operatorname{dom}(\\varphi_F)$.\nOtherwise let $\\mathbb{A}_n([G_l,v])$ to be any element in $b$.\nIt is easy to verify that this defines a randomized local algorithm.\n\nLet $G\\in \\operatorname{\\bf FG}^{S}$ be such that $|V(G)|=n$ and\n$$[\\mathcal{B}_{G}(v,\\ell_n+t),v]=\\left[\\mathcal{B}({\\bf 0}_d,\\ell_n+t),{\\bf 0}_d\\right]$$\nholds for every $v\\in V(G)$.\nNote that by the assumption on $g$ such graphs do exist, e.g., large enough torus.\nFor every $v\\in V(G)$ define the event\n$$A_v=\\left\\{\\ell\\in [0,1]^{V(G)}:\\ell\\upharpoonright \\mathcal{B}_G(v,\\ell_n)\\not \\in \\operatorname{dom}(\\varphi_F)\\right\\}.$$\nNote that if $\\ell\\not \\in \\bigcup_{v\\in V(G)} A_v$, then $\\mathbb{A}(G_l)$ is a $\\Pi$-coloring.\nTo see this pick $v\\in V(G)$ and any extension $x\\in [0,1]^{\\mathbb{Z}^d}$ of $\\ell\\upharpoonright \\mathcal{B}_G(v,\\ell_n+t)$ that is in the domain of $F$, where we interpret $v$ as the origin ${\\bf 0}_d$.\nThen $F(x)\\in X_\\Pi$ and\n$$F(x)\\upharpoonright \\mathcal{B}({\\bf 0}_d,t)\\simeq\\mathbb{A}(G_l)\\upharpoonright \\mathcal{B}_G(v,t).$$\nConsequently, we have $\\mathcal{P}([\\mathcal{B}_{G_{\\mathbb{A}(G_l)}}(v,t),v])=1$.\n\nLet $v\\in V(G)$.\nWe have $\\ell\\in A_v$ if and only if $R_F(x)>\\ell_n$ for every $x\\in [0,1]^\\Gamma$ that extends $\\ell$, i.e., every $x\\in [0,1]^{\\mathbb{Z}^d}$ that satisfies\n$$[B_{G_l}(v,\\ell_n),v]=[B_{\\mathbb{Z}^d_x}({\\bf 0}_d,\\ell_n),{\\bf 0}_d].$$ \nThis implies that\n$$\\lambda^{\\oplus n}(A_v)=\\mu\\left(\\left\\{x\\in [0,1]^{\\mathbb{Z}^d}:R_F(x)>\\ell_n\\right\\}\\right)\\le \\frac{1}{f(\\ell_n)}\\le \\frac{1}{n^2}.$$\nA union bound then shows that $\\lambda^{\\oplus n}\\left(\\bigcup_{v\\in V(G)}A_v\\right)\\le \\frac{1}{n}$ and the proof is finished.\n\\end{proof}\n\n\n\n\\begin{corollary}\\label{cor:tower=local}\n$\\mathsf{TAIL}(\\tow(\\Omega(r)))=\\mathsf{LOCAL}(O(\\log^* n))$.\n\\end{corollary}\n\\begin{proof}\nIt follows from Corollary~\\ref{cor:locallocal} itself has the following implications: \n\n\\begin{corollary}\n\\label{cor:block_factor}\n$\\Pi$ can be solved by a block factor (recall \\cref{def:tail}) if and only if $\\Pi \\in \\mathsf{LOCAL}(O(1))$. \n\\end{corollary}\n\\begin{proof}\nThis follows from \\cref{thm:classification_of_local_problems_on_grids,rem:chang_pettie_decomposition}. \n\\end{proof}\n\n\\begin{corollary}\n\\label{cor:BelowTOWERtrivial}\nIf $\\Pi \\in \\mathsf{TAIL}(\\mathrm{tower}(\\omega(r)))$, then $\\Pi \\in \\mathsf{LOCAL}(O(1))$ and if $\\Sigma_{in} = \\emptyset$ then $\\Pi$ is trivial. \n\\end{corollary}\n\\begin{proof}\n\\cref{lem:tail->local} gives that $\\Pi \\in\\mathsf{LOCAL}(o(\\log^* n))$ and \\cref{thm:classification_of_local_problems_on_grids} then implies the rest. \n\\end{proof}\n\nIn particular, Corollary~\\ref{cor:BelowTOWERtrivial} answers in negative Problem~(ii) in \\cite{HolroydSchrammWilson2017FinitaryColoring} that asked about the existence of a nontrivial local problem (without input labels) with super-tower function tail. \n\nThe second corollary is the following:\n\n\\begin{corollary}\n\\label{cor:local=tower}\n$\\Pi \\in \\mathsf{TAIL}(\\tow(\\Omega(r)))$ if and only if $\\Pi \\in \\mathsf{LOCAL}(O(\\log^* n))$. \n\\end{corollary}\n\\begin{proof}\nIf $\\Pi \\in \\mathsf{TAIL}(\\tow(\\Omega(r)))$, \\cref{lem:tail->local} yields $\\Pi \\in \\mathsf{RLOCAL}(O(\\log^* n))$. The other direction (special case of \\cref{thm:tail=uniform}) follows from \\cref{rem:chang_pettie_decomposition} that essentially states that it suffices to solve the independent set problem in the power graph. But this follows from a version of Linial's algorithm from \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n\\end{proof}\n\nFinally, we get the following. \n\\begin{corollary}\n\\label{cor:Below_d_local}\nIf $\\Pi \\in \\mathsf{TAIL}(\\omega(r^{2d}))$, then $\\Pi \\in \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$. \n\\end{corollary}\n\\begin{proof}\n\\cref{lem:tail->local} gives that $\\Pi \\in \\mathsf{LOCAL}(o(n^{1\/d}))$ and \\cref{thm:classification_of_local_problems_on_grids} then implies that $\\Pi \\in \\mathsf{LOCAL}(O(\\log^* n))$ which is equivalent to $\\Pi \\in \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$ by \\cref{cor:local=tower}. \n\\end{proof}\n\nCorollary~\\ref{cor:Below_d_local} answers in negative Problem~(iii) in \\cite{HolroydSchrammWilson2017FinitaryColoring} that asked about the existence of a problem with all its moments finite but with sub-tower function decay. Recall that \\cref{fact:moment_implies_tail} implies that if the $2d$-th moment is finite, then the tail is $\\omega(r^{2d})$ and, hence, the problem can in fact be solved with a tower tail decay. \nHowever, the $2d$ is not the right dependency and in the next section we prove a tight version of this result.\n\\subsection{Tight Grid Speedup for $\\mathsf{TAIL}$}\n\\label{subsec:tail_d-1_speedup}\n\nWhile \\cref{cor:Below_d_local} already answers the Problem of \\cite{HolroydSchrammWilson2017FinitaryColoring}, we now improve the speedup from the tail of $\\omega(r^{d-1})$. This result is tight (up to $o(1)$ in the exponent, see the discussion in \\cref{subsec:toast_construction} and \\cref{cor:toast_lower_bound}) and it is the first nontrivial difference between the $\\mathsf{TAIL}$ and $\\mathsf{LOCAL}$ model that we encountered so far. \n\n\n\n\n\\begin{theorem}\n\\label{thm:grid_speedup}\nLet $d,t\\in \\mathbb{N}$.\nFor $d \\ge 2$ there is $C_{d,t}>0$ such that if an LCL $\\Pi$ is in $\\mathsf{TAIL}(C_{d,t}r^{d-1})$, then $\\Pi \\in \\mathsf{TAIL}(\\tow(\\Omega(r)))$.\n\nSimilarly, for $d=1$ we have that a local coloring problem $\\Pi$ admits a solution by an oblivious random algorithm (it is in $\\mathsf{TAIL}(f(r))$ for any function $f = \\omega(1)$) if and only if $\\Pi$ is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$.\n\\end{theorem}\nNote that \\cref{fact:moment_implies_tail} hence shows that ${\\bf M}(\\Pi) > d-1$ already implies that $\\Pi$ is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$. \n\nThe proof is a simple adaptation of Lov\u00e1sz local lemma speedup results such as \\cite{chang2017time_hierarchy,fischer2017sublogarithmic}. \nThe proof goes as follows. Split the grid into boxes of side length roughly $r_0$ for some constant $r_0$. This can be done with $O(\\log^* (n))$ local complexity (or, equivalently, with tower tail decay). \n\nNext, consider a box $B$ of side length $r_0$ and volume $r_0^d$. In classical speedup argument we would like to argue that we can deterministically set the randomness in nodes of $B$ and its immediate neighborhood so that we can simulate the original oblivious local algorithm $\\mathcal{A}$ with that randomness and it finishes in at most $r_0$ steps. \nIn fact, ``all nodes of $B$ finish after at most $r_0$ steps'' is an event that happens with constant probability if the tail decay of $\\mathcal{A}$ is $r^d$, since then each node fails with probability $1\/r_0^d$ and there are $r_0^d$ nodes in $B$. \n\nHowever, we observe that it is enough to consider only the boundary of $B$ with $O(r_0^{d-1})$ nodes and argue as above only for the boundary nodes, hence the speedup from $r^{d-1}$ tail decay. \nThe reason for this is that although the coding radius of $\\mathcal{A}$ inside $B$ may be very large, just the fact that there \\emph{exists} a way of filling in the solution to the inside of $B$, given the solution on its boundary, is enough for our purposes. The new algorithm simply simulates $\\mathcal{A}$ only on the boundaries and then fills in the rest. \n\n\\begin{proof}\n\n\nWe will prove the result for $d \\ge 2$, the case $d=1$ is completely analogous and, in fact, simpler. \nSet $C=C_{d,t}:=d\\cdot t\\cdot(2\\Delta)^{2\\Delta}$, where $\\Delta=20^d$.\nSuppose that there is an oblivious local algorithm $\\mathcal{A}$ that explores the neighborhood of a given vertex $u$ and outputs a color of $u$ after seeing a radius $R$ such that \\begin{equation}\\label{eq:tail}\n\\P(R > r)\\le \\frac{1}{C r^{d-1}}\n\\end{equation}\nfor $r\\in \\mathbb{N}$ large enough.\nLet $r_0\\in \\mathbb{N}$ be large enough such that \\eqref{eq:tail} is satisfied.\nWrite $\\Box_{r_0}$ for the local coloring problem of tiling the input graph (which is either a large enough torus or an infinite grid $\\mathbb{Z}^d$) into boxes with length sizes in the set $\\{r_0,r_0+1\\}$ such that, additionally, each box $B$ is colored with a color from $[\\Delta+1]$ which is unique among all boxes $B'$ with distance at most $3(r_0+t)$. \n\nThe problem of tiling with boxes of side lengths in $\\{r_0, r_0+1\\}$ can be solved in $O(\\log^* n)$ local complexity \\cite[Theorem~3.1]{GaoJackson} but we only sketch the reason here: Find a maximal independent set of a large enough power of the input graph $\\mathbb{Z}^d$. Then, each node $u$ in the MIS collects a Voronoi cell of all non-MIS nodes such that $u$ is their closest MIS node. As the shape of the Voronoi cell of $u$ depends only on other MIS vertices in a constant distance from $u$, and their number can be bounded as a function of $d$ (see the ``volume'' argument given next), one can discretize each Voronoi cell into a number of large enough rectangular boxes. Finally, every rectangular box with all side lengths of length at least $r^2$ can be decomposed into rectangular boxes with all side lengths either $r$ or $r+1$. \n\nNote that in additional $O(\\log^* n)$ local steps, the boxes can then be colored with $(\\Delta+1)$-many colors in such a way that boxes with distance at most $3(r_0+t)$ must have different colors.\nTo see this, we need to argue that there are at most $\\Delta$ many boxes at distance at most $3(r_0+t)$ from a given box $B$.\nThis holds because if $B,B'$ are two boxes and $d(B,B')\\le 3(r_0+t)$, then $B'$ is contained in the cube centered around the center of $B$ of side length $10(r_0+t)$.\nThen a volume argument shows that there are at most $\\Delta\\ge \\frac{(10(r_0+t))^d}{r_0^d}$ many such boxes $B'$ (when $B$ is fixed).\n\nWe show now that there is a local algorithm $\\mathcal{A'}$ that solves $\\Pi$ when given the input graph $G$ together with a solution of $\\Box_{r_0}$ as an input.\nThis clearly implies that $\\Pi$ is in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\n\n\n\nLet $\\{B_i\\}_{i\\in I}$ be the input solution of $\\Box_{r_0}$, together with the coloring $B_i\\mapsto c(B_i)\\in [\\Delta+1]$.\nWrite $V(B)$ for the vertices of $B$ and $\\partial B$ for the vertices of $B$ that have distance at most $t$ from the boundary of $B$.\nWe have $|\\partial B|< d\\cdot t \\cdot (r_0+1)^{d-1}$.\nDefine $N(B)$ to be the set of rectangles of distance at most $(r_0+t)$ from $B$.\n\n\\begin{figure}\n \\centering\n \\includegraphics[width= .9\\textwidth]{fig\/tail_LLL.eps}\n \\caption{The picture shows one block $B$ with a red boundary $\\partial B$ of length $t$. Boxes in $N(B)$, that is with distance at most $(r_0+t)$ from $B$ are dark grey. }\n \\label{fig:tail_LLL}\n\\end{figure}\n\nWe are now going to describe how $\\mathcal{A}'$ iterates over the boxes in the order of their colors and for each node $u$ in the box $B$ it carefully fixes a random string $\\mathfrak{r}(u)$. \nThis is done in such a way that after fixing the randomness in all nodes, we can apply the original algorithm $\\mathcal{A}$ to each vertex in $\\partial B$ for each $B$ with the randomness given by $\\mathfrak{r}$ and the local algorithm at $u \\in \\partial B$ will not need to see more than its $r$-hop neighborhood. \n\nGiven a partial function $\\mathfrak{r_\\ell}:V(G)\\to [0,1]$ corresponding to nodes such that their randomness was already fixed, every box $B$ can compute a probability $\\mathbb{P}(\\fE_{\\ell}(B))$, where $\\fE_{\\ell}(B)$ is the (bad) event that there is a vertex in $\\partial B$ that needs to look further than $r_0$ to apply the local algorithm $\\mathcal{A}$.\nNote that for computing $\\mathbb{P}(\\fE_{\\ell}(B))$ it is enough to know the $r_0$-neighborhood of $B$.\n\nThe algorithm $\\mathcal{A}$ has $\\Delta+1$ rounds.\nIn the $0$-th round, we set $\\mathfrak{r}_0$ to be the empty function.\nNote that the probability that a vertex needs to look further than $r_0$ in the algorithm $\\mathcal{A}$ is smaller than $\\frac{1}{C r_0^{d-1}}$ by \\eqref{eq:tail}.\nThe union bound gives $\\mathbb{P}(\\fE_0(B))<\\frac{1}{(2\\Delta)^{2\\Delta}}$ for every $B$.\n\nIn each step $\\ell\\in [\\Delta+1]$ we define $\\mathfrak{r}_B:V(B)\\to [0,1]$ for every $B$ such that $c(B)=\\ell$ and put $\\mathfrak{r}_l$ to be the union of all $\\mathfrak{r}_D$ that have been defined so far, i.e., for $D$ such that $c(D)\\le \\ell$.\nWe make sure that the following is satisfied for every $B$:\n\\begin{enumerate}\n\t\\item [(a)] the assignment $\\mathfrak{r}_B$ depends only on $2(r_0+t)$-neighborhood of $B$ (together with previously defined $\\mathfrak{r}_{\\ell-1}$ on that neighborhood) whenever $c(B)=\\ell$,\n\t\\item [(b)] $\\mathbb{P}(\\fE_l(B))<\\frac{1}{(2\\Delta)^{2\\Delta-\\ell}}$,\n\t\\item [(c)] restriction of $\\mathfrak{r}_l$ to $N(B)$ can be extend to the whole $\\mathbb{Z}^d$ such that $\\mathcal{A}$ is defined at every vertex, where we view $N(B)$ embedded into $\\mathbb{Z}^d$.\n\\end{enumerate}\nWe show that this gives the required algorithm.\nIt is easy to see that by (a) this defines a local construction that runs in constantly many rounds and produces a function $\\mathfrak{r}:=\\mathfrak{r}_{\\Delta+1}:V(G)\\to [0,1]$.\nBy (b) every vertex in $\\partial B$ for some $B$ looks into its $r_0$ neighborhood and use $\\mathcal{A}$ to choose its color according to $\\mathfrak{r}$.\nFinally, by (c) it is possible to assign colors in $B\\setminus \\partial B$ for every $B$ to satisfy the condition of $\\Pi$ within each $B$.\nPicking minimal such solution for each $B\\setminus \\partial B$ finishes the description of the algorithm.\n\nWe show how to proceed from $0\\le \\ell<\\Delta+1$ to $\\ell+1$.\nSuppose that $\\mathfrak{r}_l$ satisfies the conditions above.\nPick $B$ such that $c(B)=\\ell+1$.\nNote that $\\mathfrak{r}_B$ that we want to find can influence rectangles $D\\in N(B)$ and these are influenced by the function $\\mathfrak{r}_l$ at distance at most $r_0$ from them.\nAlso given any $D$ there is at most one $B\\in N(D)$ such that $c(B)=\\ell+1$. \nWe make the decision $\\mathfrak{r}_B$ based on $2r_0$ neighborhood around $B$ and the restriction of $\\mathfrak{r}_l$ there, this guarantees that condition (a) is satisfied.\n\nFix $D\\in N(B)$ and pick $\\mathfrak{r}_B$ uniformly at random.\nWrite $\\mathfrak{S}_{B,D}$ for the random variable that computes the probability of $\\fE_{\\mathfrak{r}_B\\cup \\mathfrak{r}_l}(D)$.\nWe have $\\mathbb{E}(\\mathfrak{S}_{B,D})=\\mathbb{P}(\\fE_l(D))<\\frac{1}{(2\\Delta)^{2\\Delta-\\ell}}$.\nBy Markov inequality we have\n$$\\mathbb{P}\\left(\\mathfrak{S}_{B,D}\\ge 2\\Delta \\frac{1}{(2\\Delta)^{2\\Delta-\\ell}}\\right)<\\frac{1}{2\\Delta}.$$\nSince $|N(B)|\\le \\Delta$ we have that with non-zero probability $\\mathfrak{r}_B$ satisfies $\\mathfrak{S}_{B,D}<\\frac{1}{(2\\Delta)^{2\\Delta-\\ell-1}}$ for every $D\\in N(B)$.\nSince $\\mathcal{A}$ is defined almost everywhere we find $\\mathfrak{r}_B$ that satisfies (b) and (c) above.\nThis finishes the proof. \n\\end{proof}\n\nAs we already mentioned, the proof is in fact an instance of the famous Lov\u00e1sz Local Lemma problem. This problem is extensively studied in the distributed computing literature \\cite{chang2016exp_separation, fischer2017sublogarithmic, brandt_maus_uitto2019tightLLL, brandt_grunau_rozhon2020tightLLL} as well as descriptive combinatorics \\cite{KunLLL,OLegLLL,MarksLLL,Bernshteyn2021LLL}.\nOur usage corresponds to a simple LLL algorithm for the so-called exponential criteria setup (see \\cite{fischer2017sublogarithmic} and \\cite{Bernshteyn2021local=cont}). \n\n\n \n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{$\\TOAST$: General Results}\n\\label{sec:Toast}\n\nIn this section we introduce the toast construction. This construction is very standard in descriptive combinatorics \\cite{FirstToast,GJKS,Circle} and the theory of random processes (there it is often called \\emph{the hierarchy}) \\cite{HolroydSchrammWilson2017FinitaryColoring,ArankaFeriLaci,Spinka}.\n\nWhile the construction itself is well-known, in this paper we study it as a class of problems by itself. An analogy to our work is perhaps the paper of \\cite{ghaffari2018derandomizing} that defines the class $\\mathsf{S-LOCAL}$ and shows how to derandomize it.\n\n\n\nInformally, the basic toast construction works as follows. There is an adversary that has access to the (infinite) input grid $\\mathbb{Z}^d$. The adversary is sequentially choosing some connecting pieces of the grid and gives them to our algorithm $\\mathcal{A}$. The algorithm $\\mathcal{A}$ needs to label all vertices in the piece $D$ once it is given $D$. \nNote that $\\mathcal{A}$ does not know where exactly $D$ lies in $\\mathbb{Z}^d$ and its decision cannot depend on the precise position of $D$ in the input graph. In general, the piece $D$ can be a superset of some smaller pieces that were already colored by $\\mathcal{A}$. In this case, $\\mathcal{A}$ needs to extend their coloring to the whole $D$. This game is played for countable number of steps until the whole grid is colored. The algorithm solves $\\Pi$ if it never happens during the procedure that it cannot extend the labeling to $D$ in a valid manner. \n\n\n\n\n\nThe construction sketched above is what we call a spontaneous deterministic toast or just $\\TOAST$ algorithm. \nIn the next subsection we give a formal definition of $\\TOAST$ algorithms and we follow by proving relations between $\\TOAST$ algorithms and $\\mathsf{TAIL}$ complexities. \n\n\n\\subsection{Formal Definition of $\\TOAST$ classes}\n\\label{subsec:toast_definition}\nWe now give a formal definition of our four $\\TOAST$ classes. First we define the hereditary structure on which the $\\TOAST$ algorithm is interpreted. This is also called a toast (we use lowercase letters) and is parameterized by the minimum distance of boundaries of different pieces of the toast. \nThis structure also corresponds to the adversary that asks our algorithm to extend a partial solution to $\\Pi$ to bigger and bigger pieces used in the informal introduction. \n\n\\begin{definition}[$q$-toast on $\\mathbb{Z}^d$]\nFor a $q >0$, a \\emph{$q$-toast on $\\mathbb{Z}^d$} is a collection $\\mathcal{D}$ of finite connected subsets of $\\mathbb{Z}^d$ that we call \\emph{pieces} such that\n\\begin{itemize}\n \\item for every $g,h\\in \\mathbb{Z}^d$ there is a piece $D\\in\\mathcal{D}$ such that $g,h\\in D$,\n \\item if $D,E\\in \\mathcal{D}$ and $d(\\partial D,\\partial E)\\le q$, then $D=E$, where $\\partial D$ is a boundary of a set $D$.\n\\end{itemize}\n\\end{definition}\n\nGiven a $q$-toast $\\mathcal{D}$, the \\emph{rank} of its piece $D \\in \\mathcal{D}$ is the level in the toast hierarchy that $D$ belongs to in $\\mathcal{D}$. Formally, we define a \\emph{rank function $\\mathfrak{r}_\\mathcal{D}$ for the toast $\\mathcal{D}$} as $\\mathfrak{r}_\\mathcal{D}:\\mathcal{D}\\to \\mathbb{N}$, inductively as\n$$\\mathfrak{r}_\\mathcal{D}(D)=\\max\\{\\mathfrak{r}_{\\mathcal{D}}(E)+1:E\\subseteq D, \\ E\\in \\mathcal{D}\\}$$ \nfor every $D\\in \\mathcal{D}$.\n\nNext, $\\TOAST$ algorithms are described in the following language of extending functions. \n\n\\begin{definition}[$\\TOAST$ algorithm]\n\\label{def:extending_function}\nWe say that a function $\\mathcal{A}$ is a \\emph{$\\TOAST$ algorithm} for a local problem $\\Pi = (\\Sigma_{in},\\Sigma_{out}, t, \\mathcal{P})$ if \n\\begin{enumerate}\n \\item [(a)] the domain of $\\mathcal{A}$ consists of pairs $(A, c)$, where $A$ is a finite connected $\\Sigma_{in}$-labeled subset of $\\mathbb{Z}^d$ and $c$ is a partial map from $A$ to $\\Sigma_{out}$,\n \\item [(b)] $\\mathcal{A}(A,c)$ is a partial coloring from $A$ to $\\Sigma_{out}$ that extends $c$, i.e., $\\mathcal{A}(A,c)\\upharpoonright \\operatorname{dom}(c)=c$, and that depends only on the isomorphism type of $(A,c)$ (not on the position in $\\mathbb{Z}^d$). \n \n\\end{enumerate}\nA \\emph{randomized $\\TOAST$ algorithm} is moreover given as an additional input a labeling of $A$ with real numbers from $[0,1]$ (corresponding to a stream of random bits). \nFinally, if $\\mathcal{A}$ has the property that $\\operatorname{dom}(\\mathcal{A}(A,c)) =A$ for every pair $(A,c)$, then we say that $\\mathcal{A}$ is a \\emph{spontaneous (randomized) $\\TOAST$ algorithm}.\n\\end{definition}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=\\textwidth]{fig\/3-col.eps}\n \\caption{A $q$-toast and an example of an extending function for $3$-coloring. Observe that our extending function does not give a correct algorithm for $3$-coloring since the piece on the right is contained in an even bigger piece and the $3$-coloring constructed on the right cannot be extended to the node marked with an arrow}\n \\label{fig:toast_example}\n\\end{figure}\n\n\nIt is easy to see that a $\\TOAST$ algorithm $\\mathcal{A}$ can be used to color ($\\Sigma_{in}$-labeled) $\\mathbb{Z}^d$ inductively on the rank $\\mathfrak{r}_\\mathcal{D}$ of any $q$-toast $\\mathcal{D}$ for every $q\\in \\mathbb{N}$. \nNamely, define $\\mathcal{A}(\\mathcal{D})$ to be the partial coloring of $\\mathbb{Z}^d$ that is the union of $\\mathcal{A}(D)$, where $D\\in \\mathcal{D}$ and $\\mathcal{A}(D)$ is defined inductively on $\\mathfrak{r}_\\mathcal{D}$.\nFirst put $\\mathcal{A}(D):=\\mathcal{A}(D,\\emptyset)$ whenever $\\mathfrak{r}_\\mathcal{D}(D)=0$.\nHaving defined $\\mathcal{A}(D)$ for every $D\\in \\mathcal{D}$ such that $\\mathfrak{r}_\\mathcal{D}(D)=k$, we define\n$$\\mathcal{A}(D'):=\\mathcal{A}\\left(\\bigcup\\{E\\in \\mathcal{D}:E\\subseteq D'\\},\\bigcup\\{\\mathcal{A}(E):E\\subseteq D', \\ E\\in \\mathcal{D}\\}\\right).$$\nfor every $D'\\in \\mathcal{D}$ such that $\\mathfrak{r}_\\mathcal{D}(D')=k+1$. That is, we let $\\mathcal{A}$ to extend the partial solution to $\\Pi$ to bigger and bigger pieces of the toast $\\mathcal{D}$. \nObserve that if $\\mathcal{A}$ is a spontaneous $\\TOAST$ function, then $\\operatorname{dom}(\\mathcal{A}(\\mathcal{D}))=\\mathbb{Z}^d$ for every $q$-toast $\\mathcal{D}$ and every $q\\in \\mathbb{N}$.\n\n\\paragraph{$\\TOAST$ Algorithms}\nWe are now ready to define the four types of toast constructions that we are going to use (cf. \\cref{fig:big_picture_grids}). \n\n\\begin{definition}[$\\LTOAST$]\nLet $\\Pi$ be an LCL.\nWe say that $\\Pi$ is in the class $\\LTOAST$ if there is $q>0$ and an $\\TOAST$ algorithm $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is a $\\Pi$-coloring for every $q$-toast $\\mathcal{D}$.\n\nMoreover, we say that $\\Pi$ is in the class $\\TOAST$ if there is $q\\in \\mathbb{N}$ and a spontaneous $\\TOAST$ algorithm $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is a $\\Pi$-coloring for every $q$-toast $\\mathcal{D}$.\n\\end{definition}\n\nBefore we define a randomized version of these notions, we make several comments.\nFirst note that $\\TOAST\\subseteq \\LTOAST$.\nBy (c), if we construct a ($q$-)toast using an oblivious local algorithm, then we get instantly that every problem from the class $\\TOAST$ is in the class $\\mathsf{TAIL}$, see Section~\\ref{sec:Toast}.\nSimilar reasoning applies to the class $\\mathsf{BOREL}$, see Section~\\ref{sec:OtherClasses}.\nWe remark that we are not aware of any LCL that is not in the class $\\TOAST$.\n\n\\begin{comment}\n\\begin{definition}[Randomized toast algorithm]\n\\label{def:extending_function}\nLet $b$ be a set.\nWe say that a function $\\mathcal{A}$ is a randomized $\\TOAST$ algorithm if \n\\begin{enumerate}\n \\item [(a)] the domain of $\\mathcal{A}$ consists of triplets $(A, c, \\ell)$, where $A$ is a finite connected subset of $\\mathbb{Z}^d$, $c$ is a partial coloring of $A$ and $\\ell:A\\to [0,1]$. \n \\item [(b)] $\\mathcal{A}(A,c,\\ell)$ is a partial coloring of $A$ that extends $c$, i.e., $\\mathcal{A}(A,c,\\ell)\\upharpoonright \\operatorname{dom}(c)=c$. \n \\item [(c)] the function $\\mathcal{A}$ is invariant with respect to the shift action of $\\mathbb{Z}^d$ on $[0,1]^{\\mathbb{Z}^d)}$. \n\\end{enumerate}\nMoreover, if $\\mathcal{A}$ has the property that $\\operatorname{dom}(\\mathcal{A}(A,c,\\ell)) =A$ for every triplet $(A,c,\\ell)$, then we say that $\\mathcal{A}$ is a \\emph{spontaneous randomized toast algorithm}.\n\\end{definition}\n\nSimilarly as above, we define $\\mathcal{A}(\\mathcal{D},\\ell)$ for every $q$-toast $\\mathcal{D}$, $q\\in \\mathbb{N}$ and $\\ell\\in [0,1]^{\\mathbb{Z}^d}$.\nObserve that if $\\mathcal{A}$ is a spontaneous randomized toast algorithm, then $\\operatorname{dom}(\\mathcal{A}(\\mathcal{D},\\ell))=\\mathbb{Z}^d$.\n\\end{comment}\n\n\\begin{definition}[$\\RLTOAST$]\nLet $\\Pi$ be an LCL.\nWe say that $\\Pi$ is in the class $\\RLTOAST$ if there is $q>0$ and a randomized $\\TOAST$ algorithm $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is a $\\Pi$-coloring for every $q$-toast $\\mathcal{D}$ with probability $1$.\nThat is, after $\\Sigma_{in}$-labeling is fixed, with probability $1$ an assignment of real numbers from $[0,1]$ to $\\mathbb{Z}^d$ successfully produces a $\\Pi$-coloring against all possible toasts $\\mathcal{D}$ when $\\mathcal{A}$ is applied.\n\nMoreover, we say that $\\Pi$ is in the class $\\RTOAST$ if there is $q\\in \\mathbb{N}$ and a spontaneous randomized $\\TOAST$ algorithm $\\mathcal{A}$ as above.\n\\end{definition}\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{$\\mathsf{TAIL}$ Construction of $\\TOAST$}\n\\label{subsec:toast_construction_in_tail_sketch}\n\n\nThe class $\\TOAST$ (or $\\RTOAST$) is especially useful since in \\cref{sec:SpecificProblems} we discuss how many local problems that are not in class $O(\\log^* n)$ but on the other hand are not ``clearly hard'' like vertex $2$-coloring are in the class $\\TOAST$. Importantly, in \\cref{subsec:toast_construction} we show how one can construct the toast decomposition with an oblivious local algorithm in $\\mathsf{TAIL}(r^{1 - O(1\/\\sqrt{\\log r})})$, which yields the following theorem.\n\n\\begin{restatable}{theorem}{toastconstruction}\n\\label{thm:toast_in_tail}\nLet $\\Pi$ be an LCL that is in the class $\\TOAST$.\nThen \\[\n\\Pi \\in \\mathsf{TAIL}\\left(\\frac{r}{2^{O(\\sqrt{\\log r})}}\\right) = \\mathsf{TAIL}\\left(r^{1-o(1)}\\right)\n\\]\nand, hence, ${\\bf M}(\\Pi)\\ge 1$.\n\\end{restatable}\n\nWe observe in \\cref{cor:toast_lower_bound} that faster then linear tail decay is not possible. \n\n\\paragraph{Proof Sketch}\n\nHere we provide an intuition for the proof that is otherwise technical. We first check that maximal independent sets of power graphs of the grid $G = \\mathbb{Z}^d$ can be constructed effectively, e.g. because of the MIS algorithm of Ghaffari \\cite{ghaffari2016MIS}. The slowdown by this construction can be neglected in this high-level overview. \n\nWe then take a sequence $r_k = 2^{k^2}$ and construct sparser and sparser MIS in graphs $G^{r_k}, k = 1, 2, \\dots$. After constructing MIS of $G^{r_k}$, we construct Voronoi cells induced by this MIS. These Voronoi cells $V_1, V_2, \\dots$ will be, roughly, the pieces of rank $k$, but we need to ensure that all the boundaries are at least $q$ hops apart. \n\nTo do so, we first shrink each Voronoi cells by $t_k = O(\\sum_{j < k} r_j) \\ll r_k$ hops (grey boundary area in \\cref{fig:toast_construction}) and get new, smaller, cells $C_1, C_2, \\dots$. The final sequence of pieces $H_1, H_2, \\dots$ is then defined inductively, namely to define $H_i$ we start with $C_1$, then add all Voronoi cells from $k-1$-th iteration that $C_1$ intersects, then add all Voronoi cells from $k_2$-th iteration that the new, bigger, cluster intersects, \\dots\nThe definition of $t_k$ ensures that the final pieces constructed after all $k$ iterations are sufficiently far from each other, and also compatible with the smaller rank pieces. In fact, we always have $H_i \\subseteq V_i$. \n\nIt remains to argue that after $k$ steps of this process, only a small fraction of nodes remains uncovered. This is done roughly as follows. The vertices from $V_i \\setminus C_i$ that are not covered by $C_i$s are roughly $t_k \/ r_k \\approx A r_{k-1} \/ r_k$ fraction of vertices where $A$ is some (large) constant. By induction, $1 - A r_{k-2} \/ r_{k-1}$ fraction of those vertices were covered in step $k-1$ and so on. The fraction of remaining vertices is then $\\frac{A r_{k-1}}{r_{k}} \\cdot \\frac{A r_{k-2}}{r_{k-1}} \\dots \\frac{A }{r_1} = \\frac{A^k}{r_k}$. This fraction is asymptotically almost equal to $1\/r_k$, which is the desired tail behavior. However, the additional term $A^k$ is the reason why the final tail decay for $r$ such that $r_k \\le r \\le r_{k+1} = 2^{(k+1)^2}$ is bounded by $2^{O(k)}\/2^{k^2} \\le \\left( \\frac{1}{r} \\right)^{\\frac{k^2 - O(k)}{(k+1)^2}} = \\left( \\frac{1}{r} \\right)^{1 - O(1\/k)} = \\left( \\frac{1}{r} \\right)^{1 - O(1\/\\sqrt{\\log r})} $, as needed. \n\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width = .95\\textwidth]{fig\/toast_construction.eps}\n \\caption{A step during the toast construction. Voronoi cell $V$ is first shrinked to $C$ and we get the final piece $H$ by taking the union of smaller-level Voronoi cells intersecting $C$ (in this picture, there is only one lower level but the construction is recursive in general). }\n \\label{fig:toast_construction}\n\\end{figure}\n\n\n\n\n\n\n\n\n\n\n\\subsection{$\\RTOAST = \\TOAST$}\n\\label{subsec:rtoast=toast}\n\n\n\nIn this section we show that randomness does not help spontaneous toast algorithms.\n\n\\begin{restatable}{theorem}{toastrtoast}\\label{thm:derandomOfToast}\nLet $\\Pi$ be an LCL.\nThen $\\Pi$ is in the class $\\RTOAST$ if and only if $\\Pi$ is in the class $\\TOAST$.\n\\end{restatable}\n\nThe proof goes as follows. We start with a random toast algorithm $\\mathcal{A}$ and make it deterministic as follows. When the deterministic toast algorithm gets a piece $D$ of rank $1$, it considers all possible outputs that the algorithm $\\mathcal{A}$ could produce given the input random string and it uses any one of them that $\\mathcal{A}$ would produce with nonzero probability. \nBy induction, whenever $\\mathcal{A}'$ later gets a piece $D$ containing some pieces of smaller rank with a partial solution that $\\mathcal{A}'$ already produced, the original algorithm $\\mathcal{A}$ could have produced the same partial solution with nonzero probability, hence $\\mathcal{A}$ can extend the coloring in a valid manner and $\\mathcal{A}'$ deterministically choosing any such valid extension that $\\mathcal{A}$ would use with nonzero probability. \nThis shows that $\\mathcal{A}'$ also solves $\\Pi$. \nA formal proof follows. \n\n\\begin{proof}\n\nIt is easy to see that $\\TOAST\\subseteq \\RTOAST$.\nSuppose that $\\Pi$ is in the class $\\RTOAST$ and fix a witnessing $q\\in \\mathbb{N}$ together with a spontaneous randomized $\\TOAST$ algorithm $\\mathcal{A}$.\nWe may assume that $t$, the locality of $\\Pi$, satisfies $t\\le q$.\n\nLet $(A,c)$ be such that $A\\subseteq \\mathbb{Z}^d$ is finite and connected and $\\operatorname{dom}(c)\\subseteq A$.\nWe say that $(A,c)$ comes from a \\emph{partial random $q$-toast} if there is a collection $\\mathcal{D}$ of subsets of $A$ together with an event $\\mathcal{F}$ on $\\operatorname{dom}(c)$ such that\n\\begin{itemize}\n \\item if $D,E\\in \\mathcal{D}$ and $d(\\partial D,\\partial E)\\le q$, then $D=E$,\n \\item $A\\in \\mathcal{D}$,\n \\item $\\bigcup_{E\\in \\mathcal{D}\\setminus \\{A\\}}E=\\operatorname{dom}(c)$,\n \\item $\\mathcal{F}$ has non-zero probability and if we apply inductively $\\mathcal{A}$ on $\\operatorname{dom}(c)$ we obtain $c$.\n\\end{itemize}\nIf $(A,c)$ comes from a partial random $q$-toast, then we let $\\mathcal{E}(A,c)$ to be any event $\\mathcal{F}$ as above.\nOtherwise we put $\\mathcal{E}(A,c)$ to be any event on $\\operatorname{dom}(c)$ of non-zero probability.\nIt is easy to see that $\\mathcal{E}$ does not depend on the position of $A$ in $\\mathbb{Z}^d$.\nThis is because $\\mathcal{A}$ is such.\n\nWe now define a spontaneous toast algorithm $\\mathcal{A}'$.\nGiven $(A,c)$, we let $\\mathcal{F}(A,c)$ to be any event on $A$ that has non-zero probability, satisfies\n$$\\mathcal{F}(A,c)\\upharpoonright \\operatorname{dom}(c)\\subseteq \\mathcal{E}(A,c)$$\nand $\\mathcal{A}(A,c)$ gives the same outcome on $\\mathcal{F}(A,c)$.\nLet $\\mathcal{A}'(A,c)$ be equal to this outcome on $\\mathcal{F}(A,c)$.\nIt is clear that $\\mathcal{A}'$ is a spontaneous $\\TOAST$ algorithm.\nObserve that if $(A,c)$ comes from a partial random $q$-toast, then $(A,\\mathcal{A}'(A,c))$ comes from a partial random $q$-toast.\nNamely, use the same $q$-toast witness $\\mathcal{D}$ as for $(A,c)$ together with the event $\\mathcal{F}(A,c)$.\n\nIt remains to show that $\\mathcal{A}'(\\mathcal{D})$ is a $\\Pi$-coloring for every $q$-toast $\\mathcal{D}$.\nInduction on $\\mathfrak{r}_{\\mathcal{D}}$ shows that $(D,\\mathcal{A}'(D))$ comes from a partial random $q$-toast for every $D\\in \\mathcal{D}$.\nThis is easy when $\\mathfrak{r}_{\\mathcal{D}}(D)=0$ and follows from the observation above in general.\nLet $v\\in \\mathbb{Z}^d$ and pick a minimal $D\\in \\mathcal{D}$ such that there is $D'\\in \\mathcal{D}$ that satisfies $v\\in D'$ and $D'\\subseteq D$.\nSince $t\\le q$ we have that $\\mathcal{B}(v,t)\\subseteq D$.\nSince $(D,\\mathcal{A}'(D))$ comes from a partial random $q$-toast it follows that there is a non-zero event $\\mathcal{F}$ such that \n$$\\mathcal{A}'(D)\\upharpoonright \\mathcal{B}(v,t)=\\mathcal{A}(\\mathcal{D})\\upharpoonright \\mathcal{B}(v,t)$$\non $\\mathcal{F}$.\nThis shows that the constraint of $\\Pi$ at $v$ is satisfied and the proof is finished.\n\\end{proof}\n\n\n\n\n\n\n\n\\subsection{$\\mathsf{TAIL} = \\RLTOAST$}\n\\label{subsec:tail=randomlazytoast}\n\nIn this section we prove that the class $\\RLTOAST$ in fact covers the whole class $\\mathsf{TAIL} = \\bigcup_{f = \\omega(1)} \\mathsf{TAIL}(f(r))$. \n\n\\begin{restatable}{theorem}{tailrltoast}\n\\label{thm:tail=rltoast}\nLet $\\Pi$ be an LCL.\nThen $\\Pi$ is in the class $\\RLTOAST$ if and only if $\\Pi$ is in the class $\\mathsf{TAIL}$.\n\\end{restatable}\n\nWe first give an informal proof. Suppose that there is a random lazy $\\TOAST$ algorithm $\\mathcal{A}$ for $\\Pi$. We construct an oblivious algorithm with nontrivial tail behavior $\\mathcal{A}'$ for $\\Pi$ by constructing a toast via \\cref{thm:toast_in_tail} and during the construction, each node moreover simulates the algorithm $\\mathcal{A}$. The algorithm $\\mathcal{A}$ finishes after some number of steps with probability $1$, hence $\\mathcal{A}'$ finishes, too, with probability $1$ for each node. \n\nOn the other hand, suppose that there is oblivious local algorithm $\\mathcal{A}'$ for $\\Pi$. \nThen we can construct a lazy randomized toast algorithm $\\mathcal{A}$ that simply simulates $\\mathcal{A}'$ for every node. This is possible since as the pieces around a given vertex $u$ that has gets are growing bigger and bigger, the radius around $u$ that $\\mathcal{A}$ sees also grows, until it is large enough and $\\mathcal{A}$ sees the whole coding radius around $u$ and applies $\\mathcal{A}'$ then. \n\nA formal proof follows. \n\n\n\\begin{proof}\nSuppose that $\\Pi$ is in the class $\\RLTOAST$, i.e, there is $q\\in\\mathbb{N}$ together with a randomized toast algorithm $\\mathcal{A}$.\nIt follows from Theorem~\\ref{thm:toast moment} that there is an oblivious local algorithm $F$ that produces $q$-toast with probability $1$.\n\nDefine an oblivious local algorithm $\\mathcal{A}'$ as follows.\nLet $\\mathcal{B}(v,r)$ be a $\\Sigma_{in}$-labeled neighborhood of $v$.\nAssign real numbers from $[0,1]$ to $\\mathcal{B}(v,r)$ and run $F$ and $\\mathcal{A}$ on $\\mathcal{B}(v,r)$.\nLet $a_{v,r}$ be the supremum over all $\\Sigma_{in}$-labelings of $\\mathcal{B}(v,r)$ of the probability that $\\mathcal{A}'$ fails.\nIf $a_{v,r}\\to 0$ as $r \\to \\infty$, then $\\mathcal{A}'$ is an oblivious local algorithm.\n\nSuppose not, i.e., $a_{v,r}>a>0$.\nSince $\\Sigma_{in}$ is finite, we find by a compactness argument a labeling of $\\mathbb{Z}^d$ with $\\Sigma_{in}$ such that the event $\\mathcal{F}_r$ that $\\mathcal{A}'$ fails on $\\mathcal{B}(v,r)$ with the restricted labeling has probability bigger than $a$.\nNote that since the labeling is fixed, we have $\\mathcal{F}_r\\supseteq \\mathcal{F}_{r+1}$.\nSet $\\mathcal{F}=\\bigcap_{r\\in \\mathbb{N}} \\mathcal{F}_r$.\nThen by $\\sigma$-additivity we have that the probability of the event $\\mathcal{F}$ is bigger than $a>0$.\nHowever, $F$ produces toast almost surely on $\\mathcal{F}$ and $\\mathcal{A}$ produces $\\Pi$-coloring with probability $1$ for any given toast.\nThis contradicts the fact that probability of $\\mathcal{F}$ is not $0$.\n\n\n\n\\begin{comment}\nIn another words, with probability $1$ any node $v$ can reconstruct bigger and bigger piece of the toast $\\mathcal{D}$ when exploring its neighborhood.\nIn particular, $v$ will eventually became element of some piece of the toast and $\\mathcal{A}$ has to label $v$ with $\\Sigma_{out}$ after finitely many steps.\nThis defines an oblivious local algorithm $\\mathcal{B}$.\n\n\n\nPick a labeling $\\ell\\in [0,1]^{\\mathbb{Z}^d}$ such that $F(\\ell)$ is a $q$-toast and $\\mathcal{A}(\\mathcal{D},\\ell)$ is a $\\Pi$-coloring for $q$-toast $\\mathcal{D}$.\nNote that the picked labeling satisfies these conditions with probability $1$.\n\nDefine $G(\\ell)=\\mathcal{A}(F(\\ell),\\ell)$.\nIt remains to observe that $G$ is an oblivious local algorithm.\nNote that $G(\\ell)({\\bf 0}_d)=\\mathcal{A}(F(\\ell),\\ell)({\\bf 0}_d)$ depends on some $D\\in \\mathcal{D}$ such that ${\\bf 0}_d\\in D$, $\\{E\\subseteq D:E\\in \\mathcal{D}\\}$ and $\\ell\\upharpoonright D$.\nSince $F$ is an oblivious local algorithm, we see that $G$ is too.\nConsequently, $\\Pi$ is in the class $\\mathsf{TAIL}$.\n\\end{comment}\n\nThe idea to prove the converse is to ignore the $q$-toast part and wait for the coding radius.\nSuppose that $\\Pi$ is in the class $\\mathsf{TAIL}$.\nBy the definition there is an oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi$ and an unbounded function $f$ such that\n$$\\P(R_\\mathcal{A}>r)\\le \\frac{1}{f(r)}.$$\nGiven an input labeling of $\\mathbb{Z}^d$ by $\\Sigma_{in}$, we see by the definition of $R_\\mathcal{A}$ that every node $v$ must finish after finitely many steps with probability $1$.\nTherefore it is enough to wait until this finite neighborhood of $v$ becomes a subset of some piece of a given toast $\\mathcal{D}$.\nThis happens with probability $1$ independently of the given toast $\\mathcal{D}$.\n\\end{proof}\n\n\\begin{comment}\nLet $(A,c,\\ell)$ be a triplet as in definition of randomized toast algorithm.\nTake $v\\in A\\setminus \\operatorname{dom}(c)$, and suppose that $v={\\bf 0}_d$.\nOutput $\\mathcal{A}(A,c,\\ell)({\\bf 0}_d)$ if and only if there is $t\\in \\mathbb{N}$ such that $\\mathcal{B}({\\bf 0}_d,t)\\subseteq A$ and $\\ell\\upharpoonright \\mathcal{B}({\\bf 0}_d,t)$ corresponds to some coding neighborhood of $F$.\nThat is to say, any extension $\\ell\\upharpoonright \\mathcal{B}({\\bf 0}_d,t)\\sqsubseteq x\\in [0,1]^{\\mathbb{Z}^d}$ outputs the same color at $F(x)({\\bf 0}_d)$.\nIn that case output $\\mathcal{A}(A,c,\\ell)({\\bf 0}_d)=F(x)({\\bf 0}_d)$.\n\nWe show that $\\mathcal{A}$ produces a $\\Pi$-coloring for $\\mu$-almost every $\\ell\\in [0,1]^{\\mathbb{Z}^d}$ and $1$-toast $\\mathcal{D}$.\nLet $\\ell\\in [0,1]^{\\mathbb{Z}^d}$ be such that $F$ computes a $\\Pi$-coloring an need a finite coding radius for every $v\\in \\mathbb{Z}^d$.\nNote that almost every labeling $\\ell$ satisfies this condition.\nTake a $1$-toast $\\mathcal{D}$.\nIt is easy to see that $\\mathcal{A}(\\mathcal{D},\\ell)$ does not depend on $\\mathcal{D}$ and that $\\mathcal{A}(\\mathcal{D},\\ell)=F(\\ell)$, i.e., $\\mathcal{A}(\\mathcal{D},\\ell)$ is a $\\Pi$-coloring.\nThis finishes the proof.\n\\end{comment}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n \n\n\n\n\n\n\\section{$\\TOAST$: Specific Problems}\\label{sec:SpecificProblems}\n\\todo{explain that some of this is work of others and we show it to substantiate definition of toast}\n\nIn this section we focus on specific problems that are not in the class $\\mathsf{LOCAL}(O(\\log^* n))$ but they admit a solution by a (spontaneous) toast algorithm. Recall, that proving that $\\Pi \\in \\TOAST$ implies that there is an oblivious local algorithm with $1\/r^{1-o(1)}$ tail decay via \\cref{thm:toast_in_tail} and we discuss in \\cref{subsec:descriptive_combinatorics_discussion} that by a result of \\cite{GJKS2} it also implies $\\Pi\\in \\mathsf{BOREL}$. \n\nThe problems we discuss are the following: (a) vertex coloring, (b) edge coloring, (c) tiling the grid with boxes from a prescribed set (this includes, e.g., a perfect matching problem). \n\n\n\n\n\n\nAll these problems were recently studied in the literature from different points of view, see \\cite{ArankaFeriLaci, HolroydSchrammWilson2017FinitaryColoring, GJKS, spencer_personal}.\nIn all cases the main technical tool, sometimes in a slight disguise, is the toast construction.\nOur aim here is to put these problems into a unified context.\nWe note that sometimes the results in the aforementioned literature are stronger or work in a bigger generality (for example, all results discussed here can also be proven to hold on unoriented grids as discussed in \\cite{ArankaFeriLaci}, or the constructed coloring could be aperiodic \\cite{spencer_personal}, etc)\nIn what follows we only give a ``proof by picture'', sketch the description of the corresponding spontaneous toast algorithm and refer the reader to other papers for more details\/stronger results.\n\nThroughout this section we assume that $d>1$.\nWe refer the reader to \\cref{sec:OtherClasses} for definition of classes $\\mathsf{CONTINUOUS}$, $\\mathsf{BOREL}$, etc.\n\n\n\nThe main contribution if this section is the following:\n\\begin{enumerate}\n \\item The fact that the problem of vertex $3$-coloring is in $\\TOAST$ together with the speed-up \\cref{thm:grid_speedup} answers a question (i) from \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n \\item We observe that the problem of edge $2d$-coloring the grid $\\mathbb{Z}^d$ is in $\\TOAST$, thereby answering a question of \\cite{GJKS} coming from the area of descriptive combinatorics. Independently, this was proven in \\cite{ArankaFeriLaci}. \n\\end{enumerate}\n\nIn general, our approach gives that if an LCL $\\Pi$ is in the class $\\TOAST\\setminus \\mathsf{LOCAL}(O(\\log^*(n)))$, then $1\\le {\\bf M}(\\Pi)\\le d-1$ by \\cref{thm:toast_in_tail} and \\cref{thm:grid_speedup}.\nThe examples in Section~\\ref{sec:Inter} show that there are LCLs with ${\\bf M}(\\Pi)=\\ell$ for every natural number $1\\le \\ell\\le d-1$, but for most natural LCLs we do not know the exact value of ${\\bf M}(\\Pi)$, unless $d=2$.\n\n\n\\subsection{Vertex Coloring}\n\\label{subsec:3coloring}\n\nLet $d>1$.\nRecall that $\\Pi^d_{\\chi,3}$ is the \\emph{proper vertex $3$-coloring problem}.\nThis problem was studied extensively in all contexts.\nIt is known that $\\Pi^d_{\\chi,3}\\not\\in \\mathsf{LOCAL}(O(\\log^* n))$ \\cite{brandt_grids}, $\\Pi^d_{\\chi,3}\\not\\in \\mathsf{CONTINUOUS}$ \\cite{GJKS}, $\\Pi^d_{\\chi,3}\\in \\mathsf{BOREL}$ \\cite{GJKS2} and $0<{\\bf M}(\\Pi^d_{\\chi,3})\\le 2$ \\cite{HolroydSchrammWilson2017FinitaryColoring}.\nA version of the following theorem is proven in \\cite{GJKS2}, here we give a simple ``proof by picture''. \n\n\\begin{theorem}[\\cite{GJKS2}]\\label{thm:3in toast}\nLet $d>1$.\nThe LCL $\\Pi^d_{\\chi,3}$ is in the class $\\TOAST$.\n\\end{theorem}\n\n\\begin{proof}[Proof Sketch]\nThe toast algorithm tries to color all nodes in a given piece with two colors, blue and red. However, it can happen that two partial colorings of small pieces that have different parity in a bigger piece have to be connected (see \\cref{fig:3-coloring}). \nIn this case, we take the small pieces colored with incorrect parity and ``cover'' their border with the third color: green. From the inside of a small piece, the green color corresponds to one of the main two colors, while from the outside it corresponds to the other color. This enables us to switch the color parity. \n\\end{proof}\n\nAs a consequence of results from \\cref{sec:Toast} we can now fully determine the moment behavior of the problem which answers a question from \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n\n\\begin{corollary}\nLet $d>1$.\nEvery oblivious local algorithm that solves $\\Pi^d_{\\chi,3}$ has the $1$-moment infinite, but there is an oblivious local algorithm solving $\\Pi^d_{\\chi,3}$ with all $(1-\\varepsilon)$-moments finite. \nIn particular, ${\\bf M}(\\Pi^d_{\\chi,3})=1$. \n\\end{corollary}\n\\begin{proof}\nBy \\cref{thm:3in toast}, $\\Pi^d_{\\chi,3}$ is in the class $\\TOAST$.\nIt follows from \\cref{thm:toast_in_tail} that ${\\bf M}(\\Pi^d_{\\chi,3})\\ge 1$.\nBy the speedup \\cref{thm:grid_speedup} and the fact that $3$-coloring was proven not to be in class $\\mathsf{LOCAL}(O(\\log^* n)) = \\mathsf{TAIL}(\\tow(\\Omega(r))) = \\mathsf{CONTINUOUS}$ \\cite{brandt_grids,HolroydSchrammWilson2017FinitaryColoring,GJKS}, we have that every oblivious local algorithm $\\mathcal{A}$ solving $\\Pi^2_{\\chi,3}$ must have the $1$-moment infinite.\nIn particular, ${\\bf M}(\\Pi^2_{\\chi,3})=1$.\nIt remains to observe that every vertex $3$-coloring of $\\mathbb{Z}^d$ yields a vertex $3$-coloring of $\\mathbb{Z}^2$ by forgetting the last $(d-2)$ dimensions.\nConsequently, every oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi^d_{\\chi,3}$ defines an oblivious local algorithm $\\mathcal{A}'$ that solves $\\Pi^2_{\\chi,3}$ with the same tail behavior.\nThis finishes the proof.\n\\end{proof}\n\nThe $3$-coloring discussion above also implies that the toast itself cannot be constructed faster then with linear tail, otherwise $3$-coloring could be constructed with such a decay, too. \n\\begin{corollary}\n\\label{cor:toast_lower_bound}\nThere is a constant $C > 0$ such that a $q$-toast cannot be constructed in $\\mathsf{TAIL}(C r)$ on any $\\mathbb{Z}^d$.\nHere, by constructing a $q$-toast we mean that the coding radius of each node $u$ is such that $u$ can describe the minimal toast piece $D \\in \\mathcal{D}$ such that $u \\in D$. \n\\end{corollary}\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width = .95\\textwidth]{fig\/3-col-real.eps}\n \\caption{The $\\TOAST$ algorithm is given small pieces of a toast that are already vertex-colored. First step is to match the parity of colors with respect to the bigger piece. Then color the bigger piece.}\n \\label{fig:3-coloring}\n\\end{figure}\n\n\n\\subsection{Edge Coloring}\n\\label{subsec:edgecoloring}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width = .95\\textwidth]{fig\/4-edge-col.png}\n \\caption{The $\\TOAST$ algorithm is given small pieces of a toast that are already edge-colored. First step is to match the parity of colors with respect to the bigger piece. This is done by adding a three layered gadget in the problematic directions. After that the coloring is extended to the bigger piece.}\n \\label{fig:4-edge-coloring}\n\\end{figure}\\todo{dopsat}\nLet $d>1$.\nWe denote as $\\Pi^d_{\\chi',2d}$ the \\emph{proper edge $2d$-coloring problem}.\nObserve that a solution to $\\Pi^d_{\\chi',2d}$ yields a solution to the perfect matching problem $\\Pi^d_{pm}$ (because $\\mathbb{Z}^d$ is $2d$-regular).\nThis shows that $\\Pi^d_{\\chi',2d}$ is not in the class $\\mathsf{LOCAL}(O(\\log^* n))$, because the perfect matching problem clearly cannot be solved with a local algorithm: if it was, we could find a perfect matching in every large enough finite torus, which is clearly impossible since it can have an odd number of vertices \\cite{brandt_grids}. \n\n\n\n\\begin{theorem}[See also \\cite{ArankaFeriLaci}]\\label{thm:edge toast}\nLet $d>1$.\nThen $\\Pi^d_{\\chi',2d}$ is in the class $\\TOAST$.\n\\end{theorem}\n\\begin{proof}[Proof Sketch]\nThe spontaneous toast algorithm is illustrated for $d=2$ in \\cref{fig:4-edge-coloring}.\nThe ``interface'' used by the toast algorithm between smaller and bigger pieces is: smaller pieces pretend they are colored such that in the vertical direction violet and green edges alternate, while in the horizontal direction blue and red edges alternate. If all small pieces in a bigger one have the same parity, this interface is naturally extended to the bigger piece. Hence, we need to work out how to fix pieces with bad horizontal\/vertical parity, similarly to \\cref{thm:3in toast}. \n\n\\cref{fig:4-edge-coloring} shows how to fix the vertical parity, adding the highlighted gadget in the middle picture above and below given pieces keeps its interface the same but switches the parity. The same can be done in the horizontal direction. \n\nTo show the result for $d>2$ it is enough to realize that switching parity in one dimension can be done with a help of one extra additional dimension, while fixing the parity in the remaining dimensions.\nConsequently, this allows to iterate the previous argument at most $d$-many times to fix all the wrong parities.\n\\end{proof}\n\n\n\\begin{corollary}\nLet $d>1$.\nThen $1 \\le {\\bf M}(\\Pi^d_{\\chi',2d}) \\le d-1$. \nIn particular, ${\\bf M}(\\Pi^2_{\\chi',4})=1$.\n\\end{corollary}\n\nWe leave as an open problem to determine the exact value of ${\\bf M}(\\Pi^d_{\\chi',2d})$, where $d>2$.\n\nThe following corollary answers a question of \\cite{GJKS}. See \\cref{sec:OtherClasses} for the definition of class $\\mathsf{BOREL}$. \n\\begin{corollary}\nLet $d\\ge 1$.\nThen $\\Pi^d_{\\chi',2d}$ is in the class $\\mathsf{BOREL}$.\nIn particular, if $\\mathbb{Z}^d$ is a free Borel action and $\\mathcal{G}$ is the Borel graph generated by the standard generators of $\\mathbb{Z}^d$, then $\\chi'_\\mathcal{B}(\\mathcal{G})=2d$, i.e., the Borel edge chromatic number is equal to $2d$.\n\\end{corollary}\n\n\n\\begin{remark}[Decomposing grids into finite cycles]\nBencs, Hru\u0161kov\u00e1, T\u00f3th \\cite{ArankaFeriLaci} showed recently that there is a $\\TOAST$ algorithm that decomposes an unoriented grid into $d$-many families each consisting of finite pairwise disjoint cycles.\nHaving this structure immediately gives the existence of edge coloring with $2d$-colors as well as Schreier decoration.\nSince every $\\TOAST$ algorithms imply Borel solution on graphs that look locally like unoriented grids \\cite{AsymptoticDim}, this result implies that all these problems are in the class $\\mathsf{BOREL}$.\nWe refer the reader to \\cite{ArankaFeriLaci} for more details.\n\\end{remark}\n\n\n\n\n\n\n\n\\begin{comment}\n\\paragraph{Decomposing grids into finite cycles} \\todo{Jsem pro zkratit to na jeden odstavec}\n\nWe sketch how to use toast structure to decompose $\\mathbb{Z}^d$ into $d$-many families each consisting of finite pairwise disjoint cycles, formally, a collection $\\{\\mathcal{C}_i\\}_{i\\in [d]}$ such that $\\mathcal{C}_i$ is family of pairwise vertex disjoint cycles for every $i\\in [d]$ and $\\bigcup_{i\\in [d]} \\mathcal{C}_i=\\mathbb{Z}^d$.\nNote that finding such a decomposition is not an LCL because we do not have bound on the length of the cycles.\nOn the other hand, if we allow infinite paths, then such problem is LCL problem and we denote it as $\\Pi^d_{\\operatorname{cyc},d}$.\nMoreover, it is easy to see that since we work with oriented grids $\\Pi^d_{\\operatorname{cyc},d}$ is in the class $\\mathsf{LOCAL}(O(1))$.\n\nWe show that there is a spontaneous extension function $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is always a $\\Pi^d_{\\operatorname{cyc},d}$-coloring with the additional property that all the cycles in $\\mathcal{A}(\\mathcal{D})$ are finite for every $37$-toast $\\mathcal{D}$.\nMoreover, it is easy to see that our toastable algorithm also works for non-oriented grids, where the LCL problem $\\Pi^d_{\\operatorname{cyc},d}$ is not trivial.\nThe combinatorial construction is inspired by result of Bencs, Hruskova and Toth \\cite{ArankaFeriLaci} who originally proved similar result for $d=2$.\nAfter reading their paper we realized that the argument goes through for any dimension $d>1$ and we refer the reader to their paper for a detailed proof of this fact.\nHere we only sketch the main idea.\n\n\\begin{theorem}\nLet $d>1$.\nThen there is a spontaneous extension function $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is always a $\\Pi^d_{\\operatorname{cyc},d}$-coloring for every $37$-toast $\\mathcal{D}$ with the additional property that all the cycles in $\\mathcal{A}(\\mathcal{D})$ are finite.\n\\end{theorem}\n\n\\begin{remark}\nA straightforward use of the toastable algorithm together with the fact that Borel graphs with finite Borel asymptotic dimension admit Borel toast structure (see~\\cite{AsymptoticDim}) gives a decomposition of any Borel graph that looks locally like a $d$-dimensional grid (oriented or non-oriented) into $d$-many Borel families of edges each consisting of pairwise disjoint cycles.\nSince each such cycle is necessarily of even length, we get immediately that such graphs admit Borel Schreier decoration, Borel balanced orientation and their Borel chromatic index is $2d$.\n\\end{remark}\n\n\\begin{proof}[Proof Sketch]\nLet $d>1$.\nIt is easy to see that decomposing $\\mathbb{Z}^d$ into $d$ families of cycles (possibly infinite) is the same as coloring the edges with $d$-colors, such that every vertex $v\\in \\mathbb{Z}^d$ has exactly two adjacent edges of each color.\nLet $\\{\\alpha_1,\\dots,\\alpha_d\\}$ be the set of the colors.\n\nAn \\emph{(edge) coloring} of a connected finite set $A\\subseteq \\mathbb{Z}^d$ is a coloring of all edges that are adjacent to $A$.\nWe say that an edge $e$ is a \\emph{boundary edge}, if only one vertex of $e$ is contained in $A$.\nA coloring $c$ is \\emph{suitable}, if after choosing a reference point and orientation of $A$ we have that parallel boundary edges in the $i$-th direction are assigned the color $\\alpha_i$.\n\nThere is a constant $q\\in \\mathbb{N}$ such that for every $i\\not=j\\in [d]$, finite connected set $A\\subseteq \\mathbb{Z}^d$ and a suitable edge coloring $c$, there is a suitable edge coloring $c'$ of $A':=\\mathcal{B}(A,q)$ that extends $c$ such that \n\\begin{itemize}\n \\item every $c'$-cycle of color $\\alpha_i$ that contains a vertex from $A$ is finite for every $i\\in [d]$,\n \\item the $i$-th and $j$-th directions of $A$ and $A'$ are swapped, e.g., parallel $A'$-boundary edges in the $i$-th direction from the point of view of $A$ are $c'$-colored with $\\alpha_j$,\n \\item other directions remain the same.\n\\end{itemize}\nIt is easy to see that iterating this extension possibly $2d$-many times produces a spontaneous toast algorithm $\\mathcal{A}$ with all the desired properties.\n\nIt remains to sketch how to find such extension $c'$.\nFirst, since $c$ is suitable and we want to swap $i$-th and $j$-th direction, we may assume that $d=2$.\nSecond, we may assume that $A$ is a union of rectangles, such that each side has even number of vertices adjacent to a perpendicular boundary edge.\nThis could be done as follows, pick a reference point and consider decomposition of $\\operatorname{Cay}(\\mathbb{Z}^2)$ into a square of a fixed even number of vertices on a side.\nTake the union of all such squares that intersect $A$.\nThird, see Figure~\\ref{fig:SwapingColors} how to swap colors and close all cycles coming out of $A$.\nThis finishes the sketch.\n\\end{proof}\n\\end{comment}\n\n\n\n\n\\subsection{Tiling with Rectangles}\n\\label{subsec:rectangle_tiling}\n\n\nLet $d>1$ and $\\mathcal{R}$ be a finite collection of $d$-dimensional rectangles.\nWe denote as $\\Pi^d_{\\mathcal{R}}$ the problem of \\emph{tiling $\\mathbb{Z}^d$ with elements of $\\mathcal{R}$}.\nIt is not fully understood when the LCL $\\Pi^d_{\\mathcal{R}}$ is in the class $\\mathsf{LOCAL}(O(\\log^* n))$ (or equivalently $\\mathsf{CONTINUOUS}$).\nCurrently the only known condition implying that $\\Pi^d_{\\mathcal{R}}$ is \\emph{not} in $\\mathsf{LOCAL}(O(\\log^* n))$ is when GCD of the volumes of elements of $\\mathcal{R}$ is not $1$.\nMinimal unknown example is when $\\mathcal{R}=\\{2\\times 2, 2\\times 3,3\\times 3\\}$. \nFor the other complexity classes the situation is different.\nGiven a collection $\\mathcal{R}$ we define a condition \\eqref{eq:GCD} as\n\\begin{equation}\n\\label{eq:GCD}\\tag{$\\dagger$}\n \\text{GCD of side lengths in the direction of} \\ i \\text{-th generator is} \\ 1 \\ \\text{for every} \\ i\\in [d].\n\\end{equation}\nIt follows from the same argument as for vertex $2$-coloring, that if $\\mathcal{R}$ does not satisfy \\eqref{eq:GCD}, then $\\Pi^d_{\\mathcal{R}}$ is a global problem, i.e., it is not in any class from \\cref{fig:big_picture_grids}.\nOn the other hand, the condition \\eqref{eq:GCD} is sufficient for a spontaneous toast algorithm. For example, this implies $\\Pi^2_{\\{2\\times 3,3\\times 2\\}}\\in \\mathsf{BOREL} \\setminus \\mathsf{CONTINUOUS}$.\nA version of the following theorem was announced in \\cite{spencer_personal}, here we give a simple ``proof by picture''.\n\n\\begin{theorem}[\\cite{spencer_personal, stephen}]\n\\label{thm:tilings}\nLet $d>1$ and $\\mathcal{R}$ be a finite collection of $d$-dimensional rectangles that satisfy \\cref{eq:GCD}.\nThen $\\Pi^d_{\\mathcal{R}}$ is in the class $\\TOAST$.\nConsequently, ${\\bf M}(\\Pi^d_{\\mathcal{R}})\\ge 1$.\n\\end{theorem}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=.9\\textwidth]{fig\/rectangle_filling.eps}\n \\caption{The $\\TOAST$ algorithm is given small pieces of a toast that are already tiled. First step is to match the parity of the tiling with respect to the bigger piece. Then tile the bigger piece.}\n \\label{fig:Tiling}\n\\end{figure}\n\n\\begin{proof}[Proof Sketch]\nThe two-dimensional version of the construction is shown in \\cref{fig:Tiling}. The interface of pieces this time is that they cover themselves with boxes of side length $L$ which is the least common multiple of all the lengths of basic boxes. As in the previous constructions, we just need to work out how to ``shift'' each piece so that all small pieces agree on the same tiling with boxes of side length $L$ which can then be extended to the whole piece. This can be done as in \\cref{fig:Tiling}. \n\\end{proof}\n\n\n\\begin{corollary}\nLet $d=2$ and $\\mathcal{R}$ be a finite collection of $2$-dimensional rectangles.\nThen $\\Pi^2_{\\mathcal{R}}$ is either in the class $\\mathsf{LOCAL}(O(\\log^* n))$, satisfies ${\\bf M}(\\Pi^2_{\\mathcal{R}})=1$, or is global (that is, in no class in \\cref{fig:big_picture_grids}).\n\\end{corollary}\n\nFor any $d$, \\cref{thm:tilings} implies that if a tiling problem is ``hard'' (in no class in \\cref{fig:big_picture_grids}), it is because of a simple reason: a projection to one dimension gives a periodic problem. It would be interesting to understand whether there are problems that are hard but their solutions do not have any periodicity in them (see \\cref{sec:open_problems}). \n\nWe do not know if there is $\\mathcal{R}$ that would satisfy $1< {\\bf M}(\\Pi^d_{\\mathcal{R}})\\le d-1$.\nLet $d\\ge \\ell>0$ and denote as $\\mathcal{R}_\\ell$ the restriction of rectangles to first $\\ell$-dimensions.\nAs in the case of proper $3$-vertex coloring problem, we have that ${\\bf M}(\\Pi^d_{\\mathcal{R}})\\le {\\bf M}(\\Pi^\\ell_{\\mathcal{R}_{\\ell}})$.\nIn particular, if, e.g., $\\Pi^d_{\\mathcal{R}_2}$ is not in the class $\\mathsf{LOCAL}(O(\\log^* n))$, then ${\\bf M}(\\Pi^d_{\\mathcal{R}})\\le 1$.\n\n\\begin{remark}\nIn the context of non-oriented grids we need to assume that $\\mathcal{R}$ is invariant under swapping coordinates. \nSimilar argument, then shows that $\\Pi^d_\\mathcal{R}\\in \\TOAST$ if and only if $\\mathcal{R}$ satisfies \\eqref{eq:GCD}.\nCombined with the fact that there is a Borel toast on such Borel graphs (see~\\cite{AsymptoticDim}), we deduce that \\eqref{eq:GCD} for $\\mathcal{R}$ is equivalent with existence of Borel tiling with rectangles from $\\mathcal{R}$.\n\\end{remark}\n\n\n\\todo{2x3 a 3x2 se daji poskladat z elek, coz odpovida nejake partition kruznice}\n\n\n\n\n\n\\paragraph{Perfect matching}\n\nLet $d>1$.\nA special case of a tiling problem is the \\emph{perfect matching problem} $\\Pi^d_{pm}$.\nSince there are arbitrarily large tori of odd size we conclude that $\\Pi^d_{pm}\\not \\in \\mathsf{LOCAL}(O(\\log^* n))$.\nConsequently, $\\Pi^d_{pm}\\not \\in \\mathsf{TAIL}(\\tow(\\Omega(r)))$ by \\cref{thm:tail=uniform}, and it follows from the Twelve Tile Theorem of Gao et al \\cite[Theorem 5.5]{GJKS} that $\\Pi^d_{pm}\\not \\in \\mathsf{CONTINUOUS}$.\nSimilarly as with the vertex $3$-coloring problem, the proof of $\\Pi^d_{pm}\\in \\mathsf{BOREL}$, where $d>1$, from \\cite{GJKS2} uses the toast construction and, in our notation, yields the following result. \n\n\\begin{theorem}[\\cite{GJKS2}]\\label{thm:perfect_matching_intoast}\nLet $d>1$.\nThe LCL $\\Pi^d_{pm}$ is in the class $\\TOAST$.\n\\end{theorem}\n\n\\begin{corollary}\nLet $d>1$.\nThen ${\\bf M}(\\Pi^d_{pm})\\ge 1$ and every ffiid $F$ that solves $\\Pi^d_{\\chi,3}$ has the $(d-1)$-moment infinite.\nIn particular, ${\\bf M}(\\Pi^2_{pm})=1$.\n\\end{corollary}\n\nWe leave as an open problem to determine the exact value of ${\\bf M}(\\Pi^d_{pm})$, where $d>2$.\nIn fact, this problem is similar to intermediate problems that we study in \\cref{sec:Inter}. \n\n\n\n\n\n\n\n\n\n\n\n\\renewcommand{\\pm}[2]{\\operatorname{PM}^{#1, #2}_\\boxtimes}\n\\renewcommand{\\pm}[1]{\\operatorname{PM}^{#1}_\\boxtimes}\n\\newcommand{\\mathbb{Z}^d_\\boxtimes}{\\mathbb{Z}^d_\\boxtimes}\n\n\\section{Intermediate Problems}\\label{sec:Inter}\n\nIn this section we construct local problems with $\\mathsf{TAIL}$ complexities $\\Theta(r^{j-o(1)})$ for $1 \\le j \\le d-1$. \nWe first define the problem $\\pm{d}$: this is just a perfect matching problem, but in the graph $\\mathbb{Z}^d_\\boxtimes$ where two different nodes $u = (u_1, \\dots, u_d)$ and $v = (v_1, \\dots, v_d)$ are connected with an edge if $\\sum_{1 \\le i \\le d} |u_i - v_i| \\le 2$. That is, we also add some ``diagonal edges'' to the grid. We do not know whether perfect matching on the original graph (Cayley graph of $\\mathbb{Z}^d$) is also an intermediate problem (that is, not $\\mathsf{LOCAL}$ but solvable faster than $\\TOAST$). \n\nWe will show that the tail complexity of $\\pm{d}$ is $\\Theta(r^{d - 1 - o(1)})$. \nAfter this is shown, the hierarchy of intermediate problems is constructed as follows. \nConsider the problem $\\pm{d,\\ell}$. This is a problem on $\\mathbb{Z}^d$ where one is asked to solve $\\pm{d}$ in every hyperplane $\\mathbb{Z}^\\ell$ given by the first $\\ell$ coordinates in the input graph. \nThe $\\mathsf{TAIL}$ complexity of $\\pm{d,\\ell}$ is clearly the same as the complexity of $\\pm{\\ell}$, that is, $\\Theta(r^{\\ell-1-o(1)})$. \nHence, we can focus just on the problem $\\pm{d}$ from now on. \n\n\\begin{claim}\\label{cl:PM is global}\nLet $2\\le d\\in \\mathbb{N}$.\nThen the local coloring problem $\\pm{d}$ is not in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\nVia \\cref{thm:grid_speedup} this is equivalent to $\\pm{d} \\not\\in \\mathsf{TAIL}(Cr^{d-1})$ for large enough $C$. \n\\end{claim}\n\nThis follows from the fact that if there is a local algorithm, then in particular there exists a solution on every large enough torus. But this is not the case if the torus has an odd number of vertices. \n\nSecond, we need to show how to solve the problem $\\pm{d}$ faster than with a $\\TOAST$ construction algorithm. \n\n\\begin{restatable}{theorem}{intermediate}\n\\label{thm:intermediate}\nThe problem $\\pm{d}$ can be solved with an oblivious local algorithm of tail complexity $r^{d-1-O(\\frac{1}{\\sqrt{\\log r}})}$. \n\\end{restatable}\n\nHere we give only a proof sketch, the formal proof is deferred to \\cref{app:intermediate_problems}. \n\n\\begin{proof}[Proof Sketch]\nWe only discuss the case $d=3$ (and the actual proof differs a bit from this exposition due to technical reasons). As usual, start by tiling the grid with boxes of side length roughly $r_0$. \nThe \\emph{skeleton} $\\mathcal{C}_0$ is the set of nodes we get as the union of the $12$ edges of each box. \nIts \\emph{corner points} $X_0 \\subseteq \\mathcal{C}_0$ is the union of the $8$ corners of each box. \n\n\nWe start by matching vertices inside each box that are not in the skeleton $\\mathcal{C}_0$. It is easy to see that we can match all inside vertices, maybe except of one node per box that we call a \\emph{pivot}. We can ensure that any pivot is a neighbor of a vertex from $\\mathcal{C}_0$. \n\nWe then proceed similarly to the toast construction of \\cref{thm:toast_in_tail} and build a sequence of sparser and sparser independent sets that induce sparser and sparser subskeletons of the original skeleton.\nMore concretely, we choose a sequence $r_k = 2^{O(k^2)}$ and in the $k$-th step we construct a new set of corner points $X_k \\subseteq X_{k-1}$ such that any two corner points are at least $r_k$ hops apart but any point of $\\mathbb{Z}^d$ has a corner point at distance at most $2r_k$ (in the metric induced by $\\mathbb{Z}^d$). \nThis set of corner points $X_k$ induces a subskeleton $\\mathcal{C}_{k} \\subseteq \\mathcal{C}_{k-1}$ by connecting nearby points of $X_k$ with shortest paths between them in $\\mathcal{C}_{k-1}$. \n\nAfter constructing $\\mathcal{C}_k$, the algorithm solves the matching problem on all nodes of $\\mathcal{C}_{k-1} \\setminus \\mathcal{C}_k$. This is done by orienting those nodes towards the closest $\\mathcal{C}_k$ node in the metric induced by the subgraph $\\mathcal{C}_k$. \nThen, the matching is found on connected components that have radius bounded by $r_k$. There is perhaps only one point in each component, a pivot, that neighbors with $\\mathcal{C}_{k}$ and remains unmatched. \n\nAfter the $k$-th step of this algorithm, only nodes in $\\mathcal{C}_k$ (and few pivots) remain unmatched. Hence, it remains to argue that only roughly a $1\/r_k^{d-1}$ fraction of nodes survives to $\\mathcal{C}_k$. This should be very intuitive -- think about the skeleton $\\mathcal{C}_k$ as roughly looking as the skeleton generated by boxes of radius $r_k$. This is in fact imprecise, because in each step the distance between corner points of $X_k$ in $\\mathcal{C}_k$ may be stretched by a constant factor with respect to their distance in the original skeleton $\\mathcal{C}_{k-1}$. This is the reason why the final volume of $\\mathcal{C}_k$ is bounded only by $2^{O(k)} \/ r_k^{d-1}$ and by a similar computation as in \\cref{thm:toast_in_tail} we compute that the tail decay of the construction is $2^{O(\\sqrt{\\log r})}\/{r^{d-1}}$. \n\\end{proof}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Other Classes of Problems}\\label{sec:OtherClasses}\n\nIn this mostly expository section we discuss other classes of problems from \\cref{fig:big_picture_grids} and how they relate to our results from previous sections. \n\n\n\n\\subsection{$\\mathsf{CONTINUOUS}$}\\label{subsec:CONT}\n\nHere we explain the fact that an LCL $\\Pi$ on $\\mathbb{Z}^d$ is in the class $\\mathsf{LOCAL}(O(\\log^*(n)))$ if and only if it is in the class $\\mathsf{CONTINUOUS}$.\nFor $\\mathbb{Z}^d$, the model $\\mathsf{CONTINUOUS}$ was extensively studied by Gao, Jackson, Krohne, Tucker-Drob, Seward and others, see \\cite{GaoJackson,GJKS,ST,ColorinSeward}.\nIn fact, the connection with local algorithms is implicit in \\cite{GJKS} under the name Twelve Tiles Theorem \\cite[Theorem~5.5]{GJKS}.\nRecently, independently Bernshteyn \\cite{Bernshteyn2021local=cont} and Seward \\cite{Seward_personal} showed that the same holds in bigger generality for all Cayley graphs of finitely generated groups.\nOur aim here is to formally introduce the model $\\mathsf{CONTINUOUS}$ and give an informal overview of the argument for $\\mathbb{Z}^d$ that follows the lines of Gao et al \\cite{GJKS}.\nThe reasons not to simply put reference to the aforementioned papers are two.\nFirst, we work with slightly more general LCLs that have inputs.\nSecond, our goal is to introduce all classes in a way accessible for a broader audience. \n\nThe setup of the $\\mathsf{CONTINUOUS}$ model is somewhat similar to the deterministic $\\mathsf{LOCAL}$ model: there is no randomness in the definition and there is a certain symmetry breaking input on the nodes. Specifically, each vertex is labelled with one label from the set $\\{0, 1\\}$; the space of all such labelings is denoted as $2^{\\mathbb{Z}^d}$. \nHowever, for the same reason as why we insist on the uniqueness of the identifiers in the $\\mathsf{LOCAL}$ model, we insist that the input $\\{0,1\\}$ labeling breaks all symmetries. Namely, we require that there is no $g$ such that shifting all elements by $g$ preserves the labeling. We write $\\operatorname{Free}(2^{\\mathbb{Z}^d}) \\subseteq 2^{\\mathbb{Z}^d}$ for this subspace. \nA \\emph{continuous algorithm} is a pair $\\mathcal{A}=(\\mathcal{C},f)$, where $\\mathcal{C}$ is a collection of $\\{0,1\\} \\times \\Sigma_{in}$-labeled neighborhoods and $f$ is a function that assigns to each element of $\\mathcal{C}$ an output label. It needs to satisfy the following. Suppose we fix a node $v$ and the input $\\{0,1\\}$-labeling $\\Lambda$ of $\\mathbb{Z}^d$ satisfies $\\Lambda \\in \\operatorname{Free}(2^{\\mathbb{Z}^d})$. If we let $v$ explore its neighborhood with an oblivious local algorithm, then we encounter an element from $\\mathcal{C}$ after finitely many steps. \nWe write $\\mathcal{A}(\\Lambda)(v)$, or simply $\\mathcal{A}(v)$, when $\\Lambda$ is understood, for the corresponding output.\n\n\\begin{comment}\n\\vspace*{5cm}\nTo break symmetries on $\\mathbb{Z}^d$ in the setting of random processes one considers iid labels on vertices.\nThe event that some symmetries remain has measure $0$ and we simply ignore such labelings.\nIn the continuous settings we label vertices with $2:=\\{0,1\\}$.\nThe space of all such labelings is denoted as $2^{\\mathbb{Z}^d}$ and we write $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ for the space of labelings that break all symmetries.\nA \\emph{continuous algorithm} is a pair $\\mathcal{A}=(\\mathcal{C},g)$, where $\\mathcal{C}$ is a collection of ${0,1},\\Sigma_{in}$-labeled neighborhoods and $g$ is a function that assigns to each element of $\\mathcal{C}$ an output label, such that if we let a node $v$ explore its labeled neighborhood $x$, then it encounters element from $\\mathcal{C}$ after finitely many steps whenever $x\\in \\Sigma_{in}^{\\mathbb{Z}^d}\\times \\operatorname{Free}(2^{\\mathbb{Z}^d})$, i.e., $x$ breaks all symmetries.\nWe write $\\mathcal{A}(x)(v)$ for the corresponding output.\n\\end{comment}\n\n\\begin{remark}\nThis is equivalent with the definition given in e.g. \\cite{GJKS} in the same way as the correspondence between oblivious local algorithms and ffiid, see \\cref{subsec:subshifts}.\nThe fact that the collection $\\mathcal{C}$ is dense, i.e., every node finishes exploring its neighborhoods in finite time, means that $\\mathcal{A}({\\bf 0})$ is \\emph{continuous} with respect to the product topology and defined in the whole space $\\operatorname{Free}(2^{\\mathbb{Z}^d})$.\nMoreover, the value computed at $v$ does not depend on the position of $v\\in \\mathbb{Z}^d$ but only on the isomorphism type of the labels, this is the same as saying that $\\mathcal{A}(\\_)$ is \\emph{equivariant}.\n\\end{remark}\n\nWe say that an LCL $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ if there is a continuous algorithm $\\mathcal{A}$ that outputs a solution for every $\\Lambda \\in \\operatorname{Free}(2^{\\mathbb{Z}^d})$.\nNote that, in general, we cannot hope for the diameters of the neighborhoods in $\\mathcal{C}$ to be uniformly bounded.\nIn fact, if that happens then $\\Pi \\in \\mathsf{LOCAL}(O(1))$. \nThis is because there are elements in $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ with arbitrarily large regions with constant $0$ label.\nApplying $\\mathcal{A}$ on that regions gives an output label that depends only on $\\Sigma_{in}$ after bounded number of steps, hence we get a local algorithm with constant local complexity.\n\nThe main idea to connect this context with distributed computing is to search for a labeling in $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ where any $\\mathcal{A}$ necessarily has a uniformly bounded diameter.\nSuch a labeling does exist!\nIn \\cite{GJKS} these labelings are called \\emph{hyperaperiodic}.\nTo show that hyperaperiodic labelings exist on $\\mathbb{Z}^d$ is not that difficult, unlike in the general case for all countable groups \\cite{ColorinSeward}. \nIn fact there is enough flexibility to encode in a given hyperaperiodic labeling some structure in such a way that the new labeling remains hyperaperiodic.\n\nWe now give a definition of a hyperaperiodic labeling. \nLet $\\Lambda \\in 2^{\\mathbb{Z}^d}$ be a labeling of $\\mathbb{Z}^d$.\nWe say that $\\Lambda$ is \\emph{hyperaperiodic} if for every $g\\in \\mathbb{Z}^d$ there is a finite set $S\\subseteq \\mathbb{Z}^d$ such that \nfor every $t\\in \\mathbb{Z}^d$ we have \n$$\\Lambda(t+S) \\not= \\Lambda(g+t+S)$$\nThis notation means that the $\\Lambda$-labeling of $S$ differs from the $\\Lambda$-labeling of $S$ shifted by $g$ everywhere. \nIn another words, given some translation $g\\in \\mathbb{Z}^d$, there is a uniform witness $S$ for the fact that $\\Lambda$ is nowhere $g$-periodic.\nIt follows from a standard compactness argument that any continuous algorithm $\\mathcal{A}$ used on a hyperaperiodic labeling finishes after number of steps that is uniformly bounded by some constant that depends only on $\\mathcal{A}$. \n\n\n\n\\begin{comment}\n\\vspace*{5cm}\n\nLet $\\Lambda \\in 2^{\\mathbb{Z}^d}$ be a labeling of $\\mathbb{Z}^d$.\nWe say that $\\Lambda$ is \\emph{hyperaperiodic} if for every $g\\in \\mathbb{Z}$ there is a finite set $S\\subseteq \\mathbb{Z}^d$ such that \nfor every $s\\in \\mathbb{Z}^d$ we have \n$$x(s+t)\\not=x(g+s+t)$$\nfor some $t\\in S$.\nIn another words, given a shift $g\\in \\mathbb{Z}^d$ there is a uniform witness $S$ for the fact that $x$ is nowhere $g$-periodic.\nIt follows from some standard compactness argument that every hyperaperiodic element needs uniformly bounded neighborhoods for each continuous algorithm $\\mathcal{A}$, the bound of course depends on the algorithm. \n\\end{comment}\n\n\\begin{theorem}[\\cite{GJKS,Bernshteyn2021local=cont}]\nLet $d>0$.\nThen an LCL $\\Pi$ on $\\mathbb{Z}^d$ is in the class $\\mathsf{CONTINUOUS}$ if and only if it is in the class $\\mathsf{LOCAL}(O(\\log^*(n)))$.\n\\end{theorem}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=.9\\textwidth]{fig\/continuous.eps}\n \\caption{On the left side we have examples of a labeled local tiling configurations compared with a possible continuity window $w$. These are then sparsely embedded into $\\mathbb{Z}^2$.}\n \\label{fig:my_label}\n\\end{figure}\n\n\\begin{proof}[Proof Sketch]\nOne direction was observed by Bernshteyn \\cite{Bernshteyn2021local=cont} and follows from the fact that there is continuous algorithm that produces a maximal $r$-independent set for every $r\\in \\mathbb{N}$ (recall \\cref{thm:classification_of_local_problems_on_grids}).\nThis is implicit, even though in a different language, already in Gao and Jackson \\cite{GaoJackson}.\n\nNext, suppose that $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ and write $\\mathcal{A}$ for the continuous algorithm for $\\Pi$. \nWe will carefully construct a certain hyperaperiodic element $\\Lambda$ and let $\\mathcal{A}$ solve $\\Pi$ on it. \nWe know that there is a certain number $w$ such that $\\mathcal{A}$ solves $\\Pi$ on $\\Lambda$ with local complexity $w$. \nWe note that even though we do not know the value $w$ in advance, it depends on $\\Lambda$ that is due to be constructed, the flexibility of the construction guarantees to ``encode'' in $\\Lambda$ all possible scenarios.\nSee the discussion at the end of the argument. \n\nThe overall approach is as follows. \nOur local algorithm $\\mathcal{A}'$ first tiles the plane with tiles of length much bigger than $w$ with $O(\\log^* n)$ local complexity. We cannot tile the plane in a periodic manner in $O(\\log^* n)$ steps, but it can be done with a constant number of different tiles such that the side length of each tile is either $r$ or $r+1$ for some $r \\gg w$ (this construction is from \\cite[Theorem~3.1]{GaoJackson}, see also the proof of \\cref{thm:grid_speedup}). \nFor each tile $T$, $\\mathcal{A}'$ labels it with certain $\\{0,1\\}$ labeling. \nThis labeling depends only on the shape of the tile.\nWe make sure that not only $T$ with its labeling, but also $T$ together with all possible local configurations of tiles around it and corresponding labelings appears in $\\Lambda$.\nIn particular, as $r\\gg t,w$ ($t$ is the radius necessary to check the solution $\\Pi$), we have that the $(w+t)$-hop neighborhood of any node is labelled in a way that also appears somewhere in the labeling $\\Lambda$.\nAfter labeling each tile, $\\mathcal{A}'$ simulates $\\mathcal{A}$. We know that $\\mathcal{A}$ finishes after $w$ steps on $\\Lambda$ and is correct there, hence $\\mathcal{A}'$ is correct and has local complexity $O(\\log^* n)$. What remains is to describe is how to get the labeling of each tile. \n\nWe use the fact that there is a lot of flexibility when constructing a hyperaperiodic labeling. \nWe define the hyperaperiodic labeling $\\Lambda$ of the input graph $\\mathbb{Z}^d$ in such a way that, after fixing $r$, each one of constantly many local configurations of $\\{r,r+1\\}$-tiles are embedded in $\\Lambda$, each one in fact in many copies that exhaust possible $\\Sigma_{in}$-labelings.\n\nThis is done as follows.\nFix any hyperaperiodic element $\\Lambda'$, these exist by \\cite[Lemma~2.8]{GJKS}.\nGiven a tile $T$, that depends on $r$, $\\{0,1\\}$-label it by any $T$-pattern that appears in $\\Lambda'$.\nConsider the collection of all possible local configurations of such labeled tiles that are moreover labeled by $\\Sigma_{in}$.\nThis collection is finite. Since we do not know $w$ in advance, do this for any value of $r$. This yields a countable collection of possible configurations. \nEmbed this collection sparsely into $\\Lambda'$, \\cref{fig:my_label}. Extend the $\\Sigma_{in}$-labeling arbitrarily. \nThis defines $\\Lambda$.\nIt is routine to check that $\\Lambda$ is hyperaperiodic, see \\cite[Lemma~5.9]{GJKS}.\n\n\n\nIt remains to recall that since $\\Lambda$ is hyperaperiodic, there is $w\\in \\mathbb{N}$ such that $\\mathcal{A}$ has uniformly bounded locality by $w$ on $\\Lambda$.\nPick $r\\gg w,t$.\nThis works as required since any $w+t$-hop neighborhood $\\mathcal{A}'$ can encounter is present in $\\Lambda$.\n\\end{proof}\n\n\\begin{comment}\n\\vspace*{4cm}\nNext, suppose that $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ and write $\\mathcal{A}$ for the continuous algorithm.\nA high-level strategy for $d=2$ is as follows.\nThere are 12 type of tiles that are parametrized by natural numbers $n,p,q\\in \\mathbb{N}$, i.e., for fixed $n,p,q$ there are $12$ tiles.\nIf $n,p,q$ are fixed, then there is a local deterministic algorithm that tiles given, large enough, torus with these tiles in such a way that the boundary parts matches.\\todo{picture here ad reference to the nice pictures in \\cite{GJKS}}\nIt is not clear at this point how large parameters we should choose.\nThe trick is to encode all such possible tiles for all $n,p,q\\in \\mathbb{N}$ in one hyperaperiodic element $x$.\nThe encoding of the tiles in $x$ guarantees that whenever the regions \\todo{picture} are the same for the fixed parameters $n,p,q$, then $\\mathcal{A}$ outputs the same label.\nSince $x$ is hyperaperiodic, we have that $\\mathcal{A}$ stops after $T$-many steps to output $\\mathcal{A}(x)(v)$ for some $T\\in \\mathbb{N}$.\nWe choose $n,p,q$ according to $T$ in such a way that for any vertex at any region there is a tile such that the $T$-neighborhood of this vertex is fully contained in the tile.\nConsequently, this defines local rule once a given, large enough, torus is tiled into these tiles of 12 types.\n\nThis description exactly follows the proof of (1) $\\rightarrow$ (2) in \\cite[Theorem~5.5]{GJKS}, where (1) is the same as saying that there is a continuous algorithm and (2) fixes the coloring of the tiles for some parameters $n,p,q$.\nHowever, we note that we are skipping some technical details about additional parameter $w$ etc.\nThe additional step to produce a local algorithm by tiling a large enough torus into these 12 tiles is standard.\nIn fact, this is proved in \\cite[Theorem~5.10]{GJKS}, again in different language, to show (2) $\\rightarrow$ (1) in \\cite[Theorem~5.5]{GJKS}.\n\nThe main difference that we encounter are the input labels.\nWe explain how to handle them and hint, simultaneously, what we mean by encoding tiles in hyperaperiodic element, again ignoring a technical parameter $w$, see \\cite[Figures~13,14,etc]{GJKS}.\nFix a hyperaperiodic element $y$.\nFor $n,p,q\\in \\mathbb{N}$ and a tile $G$, see \\eqref{fig:???}, assign to each region labels from $\\{0,1\\}$ in such a way that (a) same regions get the same pattern, (b) this pattern is part of $y$.\nConsider moreover any possible input labeling of such $\\{0,1\\}$-labeled tiles.\nNow place these labeled tiles, for every $n,p,q\\in \\mathbb{N}$, sparsely in $y$ (Figure ???) and extend the input labeling on them to the whole $\\mathbb{Z}^2$, this defines $x$.\nIt is not hard to see that $x$ remains hyperaperiodic.\n\nIn higher dimension, i.e., when $d>2$, the same argument works as well with the only difference that there are more than 12 types of tiles.\n\\end{comment}\n\n\n\n\n\n\n\n\n\n\\subsection{Descriptive Combinatorics Classes}\n\\label{subsec:descriptive_combinatorics_discussion}\n\nIn this section we define the classes of problems that are studied in descriptive combinatorics.\nThere are two equivalent ways how the notion of a \\emph{Borel graph} can be defined.\nThe less abstract one is as follows.\nConsider the unit interval $[0,1]$ with the $\\sigma$-algebra of Borel sets, that is the smallest algebra of subsets that contains sub-intervals and is closed under countable unions and complements.\nLet $A,B\\subseteq [0,1]$ be Borel sets and consider a bijection $\\varphi:A \\to B$ between them.\nWe say that $\\varphi$ is \\emph{Borel measurable} if $\\varphi^{-1}(D\\cap B)$ is a Borel set for every Borel set $D$.\nNow, let $\\mathcal{G}$ be a graph with vertex set $[0,1]$ and assume that the degree of every vertex is bounded by some $\\Delta<\\infty$.\nWe say that $\\mathcal{G}$ is a \\emph{Borel graph} if the edge relation is induced by Borel bijections, that is there are sets $A_1,\\dots, A_\\ell$, $B_1,\\dots, B_\\ell$ and Borel bijections $\\varphi_1,\\dots, \\varphi_\\ell$ such that $(x,y)$ is an edge in $\\mathcal{G}$ if and only if there is $1\\le i\\le \\ell$ such that either $\\varphi_i(x)=y$ or $\\varphi_i(y)=x$.\n\nThe more abstract way how to define Borel graphs uses the notion of a \\emph{standard Borel space}.\nFor example $[0,1]$ with the $\\sigma$-algebra of Borel sets, when we forget that it comes from the standard topology, is a standard Borel space.\nFormally, a standard Borel space is a pair $(X,\\mathcal{B})$, where $X$ is a set and $\\mathcal{B}$ is a $\\sigma$-algebra that comes from some separable and complete metric structure on $X$.\nA uniformly bounded degree graph $\\mathcal{G}$ on $(X,\\mathcal{B})$ is Borel if it satisfies the same condition as in the case of the particular space $[0,1]$.\n\nIt is easy to see that with this definition all the basic constructions from graph theory are \\emph{relatively} Borel, e.g. if $A$ is a Borel set of vertices, then $N_\\mathcal{G}(A)$ the neighbors of elements in $A$ form a Borel set.\nAlso note that since we assume that the degree is uniformly bounded, every connected component in $\\mathcal{G}$ is at most countable.\nTypical questions that are studied in this context are Borel versions of standard notions, e.g., \\emph{Borel} chromatic number, \\emph{Borel} chromatic index, \\emph{Borel} perfect matching etc.\nWe refer the reader to \\cite{pikhurko2021descriptive_comb_survey,kechris_marks2016descriptive_comb_survey} for more details.\n\nThe most relevant Borel graphs for us are the ones that are induced by a free Borel action of $\\mathbb{Z}^d$.\nThat is, there are $d$-many aperiodic Borel bijections $\\varphi_i:X\\to X$ that pairwise commute, where \\emph{aperiodic} means that $\\varphi_i^k(x)\\not=x$ for every $k>0$, $i\\in [d]$ and $x\\in X$.\nThen the Borel (labeled) graph $\\mathcal{G}$ that is induced by this action consists of edges of the form $(x,\\varphi_i(x))$, where $x\\in X$ and $i\\in [d]$.\nThe most important result in this context is that $\\mathcal{G}$ admits a \\emph{Borel toast structure}.\nThis follows from work of Gao et al \\cite{GJKS2}, see Marks and Unger \\cite{Circle} for proof.\nWe note that it is not understood in general what Borel graphs admit a Borel toast structure.\nThis is connected with the notoriously difficult question about Borel hyperfiniteness, see \\cite{AsymptoticDim}.\n\nGiven an LCL $\\Pi$, we say that $\\Pi$ is in the class $\\mathsf{BOREL}$ if there is a Borel measurable $\\Pi$-coloring of every Borel graph $\\mathcal{G}$ as above, i.e., induced by a free Borel action of $\\mathbb{Z}^d$.\nEquivalently, it is enough to decide whether $\\Pi$ can be solved on the canonical Borel graph $\\mathcal{G}$ on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$.\nMoreover, if $\\Pi$ has some inputs $\\Sigma_{in}$, then we consider any Borel labeling $X\\to \\Sigma_{in}$ and then ask if there is a Borel solution that satisfies all the constraints on $\\Sigma_{in}\\times \\Sigma_{out}$.\n\n\\begin{theorem}[Basically \\cite{GJKS}]\nLet $\\Pi$ be an LCL that is in the class $\\LTOAST$.\nThen $\\Pi\\in\\mathsf{BOREL}$.\n\\end{theorem}\n\nIn fact, all the LCL problems without input that we are aware of and that are in the class $\\mathsf{BOREL}$ are in fact in the class $\\TOAST$.\nIt is an open problem to decide whether it is always the case.\n\nAnother interesting perspective how to see Borel constructions is as a countably infinite analogue of $\\mathsf{LOCAL}(O(\\log^*(n)))$.\nNamely, suppose that you are allowed to inductively (a) construct Borel MIS with some parameter (b) do local $O(1)$-construction that depends on the previously build structure.\nIt is clear that objects that are build in this way are Borel, however, it is not clear if this uses the whole potential of Borel sets to find Borel solutions of LCLs.\n\nAs a last remark we note that traditionally in descriptive combinatorics, people consider relaxations of Borel constructions.\nThis is usually done via ignoring how a solution behave on some negligible Borel set.\nProminent examples are \\emph{measure $0$ sets} for some Borel probability measure on $(X,\\mathcal{B})$ or \\emph{meager sets} for some compatible complete metric on $(X,\\mathcal{B})$.\nWe denote the classes of problems that admit such a solution for \\emph{every} Borel probability measure or compatible complete metric as $\\mathsf{MEASURE}$ and $\\mathsf{BAIRE}$, respectively.\nWithout going into any details we mention that on general graphs these classes are much richer than $\\mathsf{BOREL}$, see e.g. \\cite{DetMarks,BrooksMeas}.\nHowever, it is still plausible that on $\\mathbb{Z}^d$ we have $\\mathsf{BOREL}=\\mathsf{BAIRE}=\\mathsf{MEASURE}$.\nIt is well-know that this is true in $\\mathbb{Z}$ \\cite{grebik_rozhon2021LCL_on_paths}.\n\n\n\\begin{comment}\nIn both cases we are allowed to neglect some small sets, where we either do not care about the solution or are allowed to use the axiom of choice.\nSuppose that $\\mathcal{G}$ is a Borel graph on a standard Borel space $(X,\\mathcal{B})$ and $\\Pi$ is an LCL.\nSuppose, in addition, that there is a Borel probability measure $\\mu$ on $(X,\\mathcal{B})$, that is a $\\sigma$-additive measure on $\\mathcal{B}$, see \\cite[Section~?]{Kec}.\nWe say that $\\Pi$ admits \n\n\n\n\n\n\n\n\n\n\nand show the arrows that are left to complete Figure~\\ref{fig:big_picture_grids}.\nRecall that a \\emph{Borel graph} $\\mathcal{G}$ is a triplet $\\mathcal{G}=(X,\\mathcal{B},E)$, where $(X,\\mathcal{B})$ is a standard Borel space, $(X,E)$ is a graph and $E\\subseteq X\\times X$ is an element of the product $\\sigma$-algebra.\nFor the purposes of this paper we always assume that $\\mathcal{G}$ is induced by a free Borel action of $\\mathbb{Z}^d$ for some $d>0$, i.e., there is a Borel action $\\mathbb{Z}^d\\curvearrowright X$ with trivial stabilizers such that $\\mathcal{G}$ is induced by the standard generators of $\\mathbb{Z}^d$.\n\nLet $\\Pi$ be an LCL.\nThen $\\Pi$ belongs to the class $\\mathsf{BOREL}$ if there is a Borel $\\Pi$-coloring of every such Borel graph $\\mathcal{G}$.\nMoreover, we say that $\\Pi$ is in the class $\\mathsf{MEASURE}$ if for every such Borel graph $\\mathcal{G}$ and for every \\emph{Borel probability measure} $\\mu$ on $(X,\\mathcal{B})$ there is a Borel $\\Pi$-coloring of $\\mathcal{G}$ on a $\\mu$-conull set.\nSimilarly, we say that $\\Pi$ is in the class $\\mathsf{BAIRE}$ if for every such Borel graph $\\mathcal{G}$ and for every \\emph{compatible Polish topology} $\\tau$ on $(X,\\mathcal{B})$ there is a Borel $\\Pi$-coloring of $\\mathcal{G}$ on a $\\tau$-comeager set.\nIt follows directly from the definition that $\\mathsf{BOREL}\\subseteq \\mathsf{MEASURE},\\mathsf{BAIRE}$.\n\\end{comment}\n\n\n\n\\begin{comment}\nNext we follow Gao et al. \\cite{GJKS}.\nDenote as $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ the free part of the Bernoulli shift $\\mathbb{Z}^d\\curvearrowright 2^{\\mathbb{Z}^d}$ endowed with the inherited product topology.\nThen it is easy to see that $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ is a Polish space and that the action is continuous.\nConsider the Borel graph $\\mathcal{G}$ on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ induced by the standard set of generators.\nThis graph is universal in the sense that an LCL $\\Pi$ is in the class $\\mathsf{BOREL}$ if and only if there is a Borel $\\Pi$-coloring of $\\mathcal{G}$, this follows from result of Seward and Tucker-Drob~\\cite{?}.\nWe say that an LCL $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ if there is a $\\Pi$-coloring of $\\mathcal{G}$ that is continuous.\nThis is the same as saying that that there is a continuous equivariant map from $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ to $X_\\Pi$.\nThe following is a corollary of the Twelve Tiles Theorem of Gao et al. \\cite{GJKS}.\n\n\\begin{theorem}\\cite[Theorem 5.5]{GJKS}\n$\\mathsf{CONTINUOUS}=\\mathsf{LOCAL}(O(\\log^* n))$.\n\\end{theorem}\n\nWe note that in the general case of all groups this equality was shown by Bernshteyn\\cite{?} and, independently, by Seward. \n\nIt remains to describe how toastable algorithms give Borel solutions.\nAs an immediate corollary we get that all the problems from Section~\\ref{sec:SpecificProblems} are in the class $\\mathsf{BOREL}$.\n\n\\begin{theorem}\n$\\LTOAST \\subseteq \\mathsf{BOREL}$.\n\\end{theorem}\n\\begin{proof}\nSuppose that $\\Pi$ is in the class $\\LTOAST$ and let $f$ be an toast algorithm with parameter $k$.\nIt follows from work of Gao et al \\cite{GJKS2}, see Marks and Unger \\cite{Circle} for proof, that every Borel graph $\\mathcal{G}$ (that is induced by a free action of $\\mathbb{Z}^d$) admits a Borel $k$-toast for every $k\\in \\mathbb{N}$.\nIt is a standard fact that using $f$ inductively on the depth of the toast produces a Borel $\\Pi$-coloring of $\\mathcal{G}$.\n\\end{proof}\n\\end{comment}\n\n\n\n\n\n\\subsection{Finitely Dependent Processes}\n\nIn this section we use the notation that was developed in \\cref{subsec:subshifts}, also we restrict our attention to LCLs without input.\nLet $\\Pi=(b,t,\\mathcal{P})$ be such an LCL.\nRecall that $X_\\Pi$ is the space of all $\\Pi$-colorings.\n\nWe say that a distribution $\\mu$ on $b^{\\mathbb{Z}^d}$ is \\emph{finitely dependent} if it is shift-invariant and there is $k\\in \\mathbb{N}$ such that the events $\\mu\\upharpoonright A$ and $\\mu \\upharpoonright B$ are independent for every finite set $A,B\\subseteq \\mathbb{Z}^d$ such that $d(A,B)\\ge k$.\nNote that shift-invariance is the same as saying that the distribution $\\mu\\upharpoonright A$ does not depend on the position of $A$ in $\\mathbb{Z}^d$ but rather on its shape.\n\nHere we study the following situation:\nSuppose that there is a finitely dependent process $\\mu$ that is concentrated on $X_\\Pi$.\nIn that case we say that $\\Pi$ is in the class $\\mathsf{FINDEP}$. \nIn other words, $\\Pi \\in \\mathsf{FINDEP}$ if there is a distribution from which we can sample solutions of $\\Pi$ such that if we restrict the distribution to any two sufficient far sets, the corresponding restrictions of the solution of $\\Pi$ that we sample are independent. \n\nIt is easy to see that all problems in $\\mathsf{LOCAL}(O(1))$ (that is, problems solvable by a block factor) are finitely dependent, but examples of finitely dependent processes outside of this class were rare and unnatural, see the discussion in \\cite{FinDepCol}, until Holroyd and Liggett \\cite{FinDepCol} provided an elegant construction for 3- and 4-coloring of a path and, more generally, for distance coloring of a grid with constantly many colors. This implies the following result.\n\n\\begin{theorem}[\\cite{FinDepCol}]\n$\\mathsf{LOCAL}(O(\\log^* n)) \\subseteq \\mathsf{FINDEP}$ (see \\cref{fig:big_picture_grids}).\n\\end{theorem} \n\nOn the other hand, recently Spinka \\cite{Spinka} proved that any finitely dependent process can be described as a finitary factor of iid process.\\footnote{His result works in a bigger generality for finitely generated amenable groups.}\nUsing \\cref{thm:toast_in_tail}, we sketch how this can be done, we keep the language of LCLs.\n\n\\begin{theorem}[\\cite{Spinka}]\nLet $\\Pi$ be an LCL on $\\mathbb{Z}^d$ and $\\mu$ be finitely dependent process on $X_\\Pi$, i.e., $\\Pi\\in \\mathsf{FINDEP}$.\nThen there is a finitary factor of iid $F$ (equivalently an oblivious local algorithm $\\mathcal{A}$) with tail decay $1\/r^{1- O(\\sqrt{\\log r})}$ such that the distribution induced by $F$ is equal to $\\mu$.\nIn fact, $\\Pi$ is in the class $\\TOAST$. \n\\end{theorem}\n\\begin{proof}[Proof Sketch]\nSuppose that every vertex is assigned uniformly two reals from $[0,1]$.\nUse the first string to produce a $k$-toast as in \\cref{thm:toast moment}, where $k$ is from the definition of $\\mu$.\n\nHaving this use the second random string to sample from $\\mu$ inductively along the constructed $k$-toast.\nIt is easy to see that this produces ffiid $F$ that induces back the distribution $\\mu$.\nIn fact, it could be stated separately that the second part describes a randomized $\\TOAST$ algorithm: always samples from $\\mu$ given the coloring on smaller pieces that was already fixed. The finitely dependent condition is exactly saying that the joint distribution over the solution on the small pieces of $\\mu$ is the same as the distribution we sampled from.\nNow use \\cref{thm:derandomOfToast} to show that $\\Pi\\in \\TOAST$. \n\\end{proof}\n\nIn our formalism, $\\mathsf{FINDEP} \\subseteq \\TOAST$, or equivalently we get $\\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r))) \\subseteq \\mathsf{FINDEP} \\subseteq \\mathsf{TAIL}(r^{1-o(1)})$.\nIt might be tempting to thing that perhaps $\\mathsf{FINDEP} =\\TOAST$.\nIt is however a result of Holroyd et al \\cite[Corollary~25]{HolroydSchrammWilson2017FinitaryColoring} that proper vertex $3$-coloring, $\\Pi^{d}_{3,\\chi}$, is not in $\\mathsf{FINDEP}$ for every $d>1$.\nThe exact relationship between these classes remain mysterious.\n\n\n\n\\begin{comment}\n\\begin{theorem}\nLet $\\Pi$ be an LCL on $\\mathbb{Z}^d$.\nSuppose that $\\Pi$ is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$.\nThen there is finitely dependent process $\\mu$ on $X_\\Pi$.\n\\end{theorem}\n\\begin{proof}\nIt follows from Theorem~\\ref{thm:classification_of_local_problems_on_grids} that there is $r,\\ell \\in \\mathbb{N}$ such that $\\ell\\le r-t$ and a rule $\\mathcal{A}$ of locality at most $\\ell$ that produces a $\\Pi$-coloring for every $(r+t)$-power proper vertex coloring of $\\mathbb{Z}^d$.\n\nSuppose that $\\mu_{(r+t)}$ is a finitely dependent $(r+t)$-power proper vertex coloring of $\\mathbb{Z}^d$ process with locality $k\\in \\mathbb{N}$.\nWe claim that composing $\\mu_{(r+t)}$ with $\\mathcal{A}$ yields a finitely dependent process with locality $(2\\ell+k)$.\nLet $A\\subseteq \\mathbb{Z}^d$ be a finite set.\nThen the output of $\\mathcal{A}$ depends on $\\mathcal{B}(A,\\ell)$.\nTherefore if $A,B\\subseteq \\mathbb{Z}^d$ satisfy $d(A,B)>2\\ell+k$, then\n$$d(\\mathcal{B}(A,\\ell),\\mathcal{B}(B,\\ell))>k$$\nand we conclude that the outputs of $\\mathcal{A}$ on $A$ and $B$ are independent.\n\nIt remains to find $\\mu_{r}$ finitely dependent $r$-power proper vertex coloring of $\\mathbb{Z}^d$ process for every $r>0$.\nThis is done by Holroyd and Liggett \\cite[Theorem~20]{FinDepCol}.\n{\\color{red} or something like that, they use a lot of colors fr $r$-power coloring...}\n\\end{proof}\n\nGiven a fiid $F$ we define $F^*\\lambda$ to be the push-forward of the product measure via $F$.\nSince $F$ is equivariant it follows that $F^*\\lambda$ is shift-invariant.\n\nThe following theorem was essentially proven by Spinka\\cite{?}. His result nicely fits to our definitions in \\cref{sec:Toast}. \n\\begin{theorem}\n\n\\end{theorem}\n\n\n\\begin{theorem}\nLet $\\Pi$ be an LCL problem on $\\mathbb{Z}^d$ and $\\mu$ be finitely dependent process on $X_\\Pi$.\nThen for every $\\epsilon>0$ there is a ffiid $F$ such that $F^*\\lambda=\\mu$ and the tail of $F$ has the $(1-\\epsilon)$-moment finite.\nMoreover, $\\Pi$ is in the class $\\TOAST$.\n\\end{theorem}\n\n\\begin{remark}\nMaybe the random toast construction works for any finitely dependent processes not just LCL.\n\nThere is the notion of strongly irreducible shifts of finite type, some comment about it...\n\nDefine $\\mathsf{FINDEP}$ to be the class of all LCL problems that admit a finitely dependent distribution $\\mu$.\nIt follows that $\\mathsf{TAIL}(\\tow(\\Omega(r)))\\subseteq \\mathsf{FINDEP} \\subseteq \\mathsf{TOAST}$ and we have $\\mathsf{FINDEP}\\not = \\TOAST$ because proper vertex $3$-coloring $\\Pi^d_{\\chi,3}$ is not in $\\mathsf{FINDEP}$ for every $d>1$, by \\cite[Corollary~25]{HolroydSchrammWilson2017FinitaryColoring}.\n\\end{remark}\n\n\n\\begin{proof}\nWe define a spontaneous randomized toast algorithm $\\Omega$ such that $\\Omega(\\mathcal{D},\\ell)$ is a $\\Pi$-coloring for every $k$-toast $\\mathcal{D}$ and every $\\ell\\in [0,1]^{\\mathbb{Z}^d}$.\nThis already implies the moreover part by Theorem~\\ref{thm:derandomOfToast}.\n\nLet $A\\subseteq \\mathbb{Z}^d$ be a finite subset of $\\mathbb{Z}^d$.\nWrite $\\mathbb{P}_A$ for the distribution of $\\Pi$-patterns of $\\mu$ on $A$, i.e., $\\mathbb{P}_A=\\mu\\upharpoonright A$.\nSuppose now that $(A,c,\\ell)$ are given, where $A\\subseteq \\mathbb{Z}^d$ is finite subset, $c$ is a partial coloring of $A$ and $\\ell\\in [0,1]^{A}$.\nConsider the conditional distribution $\\mathbb{P}_{A,c}$ of $\\mathbb{P}_A$ given the event $c$.\nIf $\\mathbb{P}_{A,c}$ is non-trivial, then define $\\Omega$ such that $\\Omega(A,c,{-})$ induces the same distribution as $\\mathbb{P}_{A,c}$ when $\\ell$ is fixed on $\\operatorname{dom}(c)$.\nOtherwise define $\\Omega$ arbitrarily.\n\nLet $\\mathcal{D}$ be a $k$-toast.\nWe show that $\\Omega(\\mathcal{D},\\ell)$ is a $\\Pi$-coloring for every $\\ell\\in [0,1]^{\\mathbb{Z}^d}$ and $\\Omega(\\mathcal{D},{-})$ produces $\\mu$.\nNote that the first claim shows that $\\Pi$ is in the class $\\RTOAST$, while the second produces $\\mu$ as ffiid as follows.\nSuppose that every $v\\in \\mathbb{Z}^d$ is assigned two independent reals $l_{v,0}$ and $l_{v,1}$.\nUse the first real to produce a $k$-toast as in Theorem~\\ref{thm:toast moment} and use the second real together with the produced toast as an input to $\\Omega$.\nIt follows from the independence of $l_{v,0}$ and $l_{v,1}$ that this produces $\\mu$.\n\nWe show that\n$$\\Omega(\\mathcal{D},{-})\\upharpoonright D=\\mu\\upharpoonright D$$\nfor every $D\\in \\mathcal{D}$ inductively on $\\mathfrak{r}_{\\mathcal{D}}$.\nIf $\\mathfrak{r}_{\\mathcal{D}}(D)=0$, then it follows from the definition.\nSuppose that $\\mathfrak{r}_{\\mathcal{D}}(D)>0$ and let $E_1,\\dots, E_t\\in \\mathcal{D}$ be an enumeration of $E\\subset D$.\nIt follows from the inductive hypothesis together with the assumption that $\\mu$ is $k$-independent and $\\mathcal{D}$ is a $k$-toast that \n$$\\Omega(\\mathcal{D},{-})\\upharpoonright \\bigcup_{i\\le t}E_i=\\mu\\upharpoonright \\bigcup_{i\\le t} E_i.$$\nNow when $\\ell$, and therefore a partial coloring $c$, is fixed on $\\bigcup_{i\\le t}E_i=\\operatorname{dom}(c)$, the definition $\\Omega(D,c,{-})$ extends $c$ in such a way that $\\mathbb{P}_{D,c}=\\mu\\upharpoonright \\operatorname{dom}(c)$.\nSince this holds for any such assignment $\\ell$ of real numbers on $\\bigcup_{i\\le t}E_i$ we deduce that\n$$\\Omega(\\mathcal{D},{-})\\upharpoonright D=\\mu\\upharpoonright D$$\nand the proof is finished.\n\\end{proof}\n\\end{comment}\n\n\n\n\n\n\n\\subsection{$\\mathsf{TAIL}$ Complexities in Distributed and Centralised Models}\n\\label{subsec:congestLCA}\n\nAlthough we are not aware of a systematic work on $\\mathsf{TAIL}$ complexities in the distributed algorithms community, similar notions like oblivious algorithms \\cite{Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring}, node-average complexity \\cite{feuilloley2017node_avg_complexity_of_col_is_const,barenboim2018avg_complexity,chatterjee2020sleeping_MIS} or energy complexity \\cite{chang2018energy_complexity_exp_separation,chang2020energy_complexity_bfs} are studied. \nThe breakthrough MIS algorithm of Ghaffari \\cite{ghaffari2016MIS} may have found many applications exactly because of the fact that its basis is an algorithm with exponential tail behavior. \n\nWe discuss two applications of algorithms with nontrivial $\\mathsf{TAIL}$ complexities. First, the node-average complexity literature studies the average coding radius that a node needs to answer. Whenever an oblivious algorithm has its first moment finite, we have an algorithm with constant average coding radius. \nWith a little more work, one can see that the result from \\cref{thm:3in toast}, that is $3$-coloring of grids has $\\mathsf{TAIL}$ complexity $r^{1-O(1\/\\sqrt{\\log r})}$ can be adapted to the setting of finite grids where it gives a $3$-coloring algorithm with average complexity $2^{O(\\sqrt{\\log n})}$. \nFor nonconstant $\\Delta$, the situation is more complicated. In \\cite{barenboim2018avg_complexity} it was shown that the Luby's $\\Delta+1$ coloring algorithm has tail $2^{\\Omega(r)}$ even for nonconstant $\\Delta$, but the question of whether MIS can have constant average coding radius in nonconstant degree graphs is open \\cite{chatterjee2020sleeping_MIS}. \n\nAnother application of the $\\mathsf{TAIL}$ bound approach comes from the popular centralised $\\mathsf{LCA}$ model \\cite{lca11,lca12}: there one has a huge graph and instead of solving some problem, e.g., coloring, on it directly, we are only required to answer queries about the color of some given nodes. It is known that, as in the $\\mathsf{LOCAL}$ model, one can solve $\\Delta+1$ coloring with $O(\\log^* n)$ queries on constant degree graphs by a variant of Linial's algorithm \\cite{even03_conflictfree}. But using the tail version of the Linial's algorithm instead, we automatically get that $O(\\log^* n)$ is only a worst-case complexity guarantee and on average we expect only $O(1)$ queries until we know the color of a given node. That is, the amortized complexity of a query is in fact constant. The same holds e.g. for the $O(\\log n)$ query algorithm for LLL from \\cite{brandt_grunau_rozhon2021LLL_in_LCA}. \n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Open Questions}\n\\label{sec:open_problems}\n\nHere we collect some open problems that we find interesting. \n\n\\begin{enumerate}\n \\item Is it true that $\\RLTOAST = \\LTOAST$? If yes, this implies $\\mathsf{TAIL} \\subseteq \\mathsf{BOREL}$. \n \\item Is it true that $\\RLTOAST = \\RTOAST$? If yes, this implies $\\mathsf{TAIL} = \\mathsf{TAIL}(r^{1-o(1)})$ and answers the fourth question of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n\t\\item Can the tail complexity $r \/ 2^{O(\\sqrt{\\log r})}$ in the construction of toast be improved to $r \/ \\operatorname{poly} \\log (r)$? The same can be asked for the complexity of intermediate problems from \\cref{sec:Inter}. \n\t\t\\item Is there a local coloring problem $\\Pi$ such that ${\\bf M}(\\Pi)$ is non-integer value or $0$, i.e., ${\\bf M}(\\Pi)\\in [0,d-1]\\setminus \\{1,2,\\dots, d-1\\}$?\n\t\t\\item Can you find an example of a local problem without input labels which does not belong to any class in \\cref{fig:big_picture_grids} but the solutions to this problem are not ``periodic'' (this cannot happen for $d=1$). Maybe (a version of) Penrose tilings can be helpful there?\n\t\\item What is the value ${\\bf M}(\\Pi^d_{pm})$, where $\\Pi^d_{pm}$ is the problem of finding a perfect matching of $\\mathbb{Z}^d$, where $d>2$?\n \\item Let $d>1$.\n\tFor what finite collections of $d$-dimensional rectangles $\\mathcal{R}$ is $\\Pi^d_{\\mathcal{R}}$ is in the class $\\mathsf{LOCAL}(O(\\log^* n))$? Is there a finite collection of $d$-dimensional rectangles $\\mathcal{R}$ that would satisfy $1< {\\bf M}(\\Pi^d_{\\mathcal{R}})\\le d-1$?\n\t\\item Is it true (on $\\mathbb{Z}^d$) that $\\mathsf{BOREL}=\\mathsf{MEASURE}$ or $\\mathsf{BOREL} = \\mathsf{BAIRE}$?\n\\end{enumerate}\n\n\\section*{Acknowledgement}\nWe would like to thank Ferenc Bencs, Anton Bernshteyn, Sebastian Brandt, Yi-Jun Chang, Mohsen Ghaffari, Aranka Hru\u0161kov\u00e1, Stephen Jackson, Oleg Pikhurko, Brandon Sewart, Yinon Spinka, L\u00e1szlo T\u00f3th, and Zolt\u00e1n Vidny\u00e1szky for many engaging discussions. \n\nThe first author was supported by Leverhulme Research Project Grant RPG-2018-424. \nThis project has received funding from the European Research Council (ERC) under the European\nUnions Horizon 2020 research and innovation programme (grant agreement No. 853109). \n\n\\bibliographystyle{alpha}\n\n\\section{Introduction}\n\nIn this paper, we study local problems on grid graphs from three different perspectives: (a) the perspective of the theory of distributed algorithms, (b) the perspective of the theory of random processes, (c) the perspective of the field of descriptive combinatorics. \nAlthough different communities in these three fields ask very similar questions, no systematic connections between them were known, until the insightful work of Bernshteyn \\cite{Bernshteyn2021LLL} who applied results from the field (a) to the field (c). In this work, our most important results follow by applying techniques from the field (a) to (b), but perhaps more importantly, we present all three different perspectives as a part of one common theory (see \\cref{fig:big_picture_grids}). \n\nIn this introductory section, we first quickly present the three fields and define the object of interest, local problems, a class of problems that contains basic problems like vertex or edge coloring, matchings, and others. \nThen, we briefly present our results, most of which can be seen as inclusions between different classes of local problems shown in \\cref{fig:big_picture_grids}. \n\n\\paragraph{Distributed Computing}\n\nThe $\\mathsf{LOCAL}$ model of computing \\cite{linial1987LOCAL} was motivated by understanding distributed algorithms in huge networks. \nAs an example, consider the network of all wifi routers, where two routers are connected if they are close enough to exchange messages. In this case, they should better communicate with user devices on different channels, so as to avoid interference. In the graph-theoretic language, we want to color the network. Even if the maximum degree is $\\Delta$ and we want to color it with $\\Delta+1$ colors, the problem remains nontrivial, because the decision of each node should be done after few communication rounds with its neighbors. \n\nThe $\\mathsf{LOCAL}$ model formalizes this setup: we have a large network, with each node knowing its size, $n$, and perhaps some other parameter like the maximum degree $\\Delta$. In case of randomized algorithms, each node has an access to a random string, while in case of deterministic algorithms, each node starts with a unique identifier from a range of size polynomial in $n$. \nIn one round, each node can exchange any message with its neighbors and can perform an arbitrary computation. We want to find a solution to a problem in as few communication rounds as possible. \nImportantly, there is an equivalent view of $t$-rounds $\\mathsf{LOCAL}$ algorithms: such an algorithm is simply a function that maps $t$-hop neighbourhoods to the final output. An algorithm is correct if and only if applying this function to each node solves the problem. \n\nThere is a rich theory of distributed algorithms and the local complexity of many problems is understood. One example: if the input graph is the $d$-dimensional torus (think of the grid $\\mathbb{Z}^d$ and identify nodes with the same remainder after dividing the coordinate with $n^{1\/d}$), a simple picture emerges. Any local problem (we define those formally later) is solvable with local complexity $O(1)$, $\\Theta(\\log^* n)$, or we need to see the whole graph to solve it, i.e., the local complexity is $\\Theta(n^{1\/d})$ (see \\cref{thm:classification_of_local_problems_on_grids}). For example, a proper coloring of nodes with $4$ colors has local complexity $\\Theta(\\log^* n)$, while coloring with $3$ colors has already complexity $\\Theta(n^{1\/d})$. \n\n\\paragraph{Finitary factors of iid processes}\n\\todo{zeptat se Spinky}\nConsider the following example of the so-called anti-ferromagnetic Ising model. \nImagine atoms arranged in a grid $\\mathbb{Z}^3$. Each atom can be in one of two states and there is a preference between neighboring atoms to be in a different state. We can formalize this by defining a potential function $V(x,y)$ indicating whether $x$ and $y$ are equal. \nThe Ising model models the behavior of atoms in the grid as follows: for the potential $V$ and inverse temperature $\\beta$, if we look at the atoms, we find them in a state $\\Lambda$ with probability proportional to $\\exp(-\\beta \\sum_{u \\sim v \\in E(\\mathbb{Z}^d)} V(\\Lambda(u), \\Lambda(v)))$. If $\\beta$ is close to zero, we sample essentially from the uniform distribution. If $\\beta$ is infinite, we sample one of the two $2$-colorings of the grid. \n\nTo understand better the behavior of phase transitions, one can ask the following question. For what values of $\\beta$ can we think of the distribution $\\mu$ over states from the Ising model as follows: there is some underlying random variable $X(u)$ at every node $u$, with these variables being independent and identically distributed. We can think of $\\mu$ as being generated by a procedure, where each node explores the random variables in its neighborhood until it decides on its state, which happens with probability $1$. If this view of the underlying distribution is possible, we say it is a finitary factor of an iid process (ffiid). \nIntuitively, being an ffiid justifies the intuition that there is a correlation decay between different sites of the grid. To understand ffiids, we try to analyze solutions to which local problems can be sampled from an ffiid distribution (2-coloring of the grid is not among them). \n\nAlthough this definition was raised in a different context then distributed computing, it is actually almost identical to the definition of a distributed algorithm! The differences are two. First, the ffiid procedure does not know the value of $n$ since, in fact, it makes most sense to consider it in infinite graphs (like the infinite grid $\\mathbb{Z}^d$). Second, the natural measure of the complexity is different: it is the speed of the decay of the random variable corresponding to the time until the local algorithm finds the solution. \nNext, note that the algorithm has to return a correct solution with probability $1$, while the usual randomized local algorithm can err with a small probability. We call ffiids \\emph{oblivious local algorithms} to keep the language of distributed computing (see \\cref{sec:preliminaries}).\nWe define the class $\\mathsf{TAIL}(f(r))$ as the class of problems solvable by an oblivious local algorithm if the probability of the needed radius being bigger than $r$ is bounded by $1\/f(r)$ (see \\cref{fig:big_picture_grids}). \n\n\\paragraph{Descriptive Combinatorics}\n\nA famous Banach-Tarski theorem states that one can split a three-dimensional unit volume ball into five pieces, translate and rotate each piece and end up with two unit volume balls. \nThis theorem serves as an introductory warning in the beginner measure theory courses: not every subset of $\\mathbb{R}^3$ can be measurable.\nSimilarly surprising was when in 1990 Laczkovich answered the famous Circle Squaring Problem of Tarski\\todo{cite laczkovich}: A square of unit area can be decomposed into finitely many pieces such that one can translate each of them to get a circle of unit area.\nUnlike in Banach-Tarski paradox, this result was improved in recent years to make the pieces measurable in various sense \\cite{OlegCircle, Circle, JordanCircle}.\n\n\nThis theorem and its subsequent strengthenings are the highlights of a field nowadays called descriptive combinatorics \\cite{kechris_marks2016descriptive_comb_survey,pikhurko2021descriptive_comb_survey} that has close connections to distributed computing as was shown in an insightful paper by Bernshteyn \\cite{Bernshteyn2021LLL}. \nGoing back to the Circle Squaring Problem, the proofs of its versions first choose a clever set of constantly many translations $\\tau_1, \\dots, \\tau_k$. Then, one constructs a bipartite graph between the vertices of the square and the cycle and connects two vertices $x, y$ if there is a translation $\\tau_i$ mapping the point $x$ in the square to the point $y$ in the cycle. The problem can now be formulated as a perfect matching problem and this is how combinatorics enters the picture. In fact, versions of a matching or flow problems are then solved on the grid graph $\\mathbb{Z}^d$ using a trick called \\emph{toast construction} also used in the area of the study of ffiids \\cite{HolroydSchrammWilson2017FinitaryColoring,ArankaFeriLaci,Spinka}. \n\nWe discuss the descriptive combinatorics classes in a more detail in the expository \\cref{sec:OtherClasses} but we will not need them in the main body. Let us now give an intuitive definition of the class $\\mathsf{BOREL}$. \nBy a result of Chang and Pettie \\cite{chang2016exp_separation}, we know that the maximal independent set problem is complete for the class $\\mathsf{LOCAL}(O(\\log^* n))$ in the sense that this problem can be solved in $O(\\log^* n)$ rounds and any problem in that class can be solved in constant number of rounds if we are given access to an oracle computing MIS (we make this precise in \\cref{rem:chang_pettie_decomposition}). Imagine now that the input graph is infinite and our local algorithm has access to MIS oracle and local computations but now both can be used countably many times instead of just once or, equivalently, fixed number of times as is the case with the class $\\mathsf{LOCAL}(O(\\log^* n))$. That is, we can think of an algorithm that constructs MIS of larger and larger powers of the input graph and uses this infinite sequence to solve an LCL. Those constructions are in the class $\\mathsf{BOREL}$ from \\cref{fig:big_picture_grids}. In fact, the exact relationship between this definition and full power of $\\mathsf{BOREL}$ is not clear.\nHowever, it should give an idea about the difference between descriptive combinatorics and distributed computing -- from the descriptive combinatorics perspective, reasonable classes are closed under repeating basic operations countably many times.\n\n\n\n\n\\paragraph{Locally Checkable Problems on Grids as a ``Rosetta Stone''}\n\nLocally checkable problems (LCLs) are graph problems where a correctness of a solution can be checked locally. That is, there is a set of local constraints such that if they are satisfied at every node, the solution to a problem is correct. \nAn examples of such constraints include vertex coloring, edge coloring, or perfect matching. We also allow inputs on each node, so the list coloring is also an example of LCL (given that the total number of colors in all lists is finite). One should think of this class as an analogue of the class NP from complexity theory. \n\nIn this paper, we study locally checkable problems on grid graphs in the three settings introduced above. This approach reveals number of connections. For example, in \\cref{fig:big_picture_grids}, we can see that the two classes of local complexities $\\mathsf{LOCAL}(O(1))$ and $\\mathsf{LOCAL}(O(\\log^* n))$ are equivalent to different natural classes from the perspective of factors of iid and descriptive combinatorics. On the other hand, those two worlds are richer. Many of our discussions in fact apply to graphs of bounded growth, but the grids serve as the simplest, but still quite nontrivial model in which one can try to understand all three settings together (the case of lines is understood quite well \\cite{grebik_rozhon2021LCL_on_paths} and we study aspects of the connection for trees in a parallel paper \\cite{brandt_chang_grebik_grunau_rozhon_vidnyaszky2021LCLs_on_trees_descriptive}). \n\n\\begin{landscape}\n\\begin{figure}\n \\centering\n \\includegraphics[width = 1.5\\textwidth]{fig\/big_picture_grids.png}\n \\caption{The big picture of local complexity classes on grids from the perspective of distributed algorithms, factors of iid, and descriptive combinatorics. Canonical example of each class is on the left side. }\n \\label{fig:big_picture_grids}\n\\end{figure}\n\\end{landscape}\n\n\n\\subsection{Our Contribution }\n\nOur contribution can be mostly summed up as several new arrows in \\cref{fig:big_picture_grids}. In general, this clarifies several basic questions regarding the connections of the three fields and helps to view all three fields as part of a general theory of locality. \n \n\n\nFor example, consider the following question from the paper of Brandt et al. \\cite{brandt_grids} that considered LCL problems on grids: \n\\begin{quote}\n\\emph{\nHowever, despite the fact that techniques seem to translate between finitary colorings and distributed complexity, it remains unclear how to directly translate results from one setting to the other in a black-box manner; for instance, can we derive the lower bound for 3-colouring from\nthe results of Holroyd et al., and does our complexity classification imply answers to the open questions they pose?\n}\n\\end{quote}\nWe show how to 'directly translate' results between the setting of distributed algorithms and finitary factors of iid in \\cref{subsec:preliminaries_tail,subsec:speedup_warmup}. This indeed answers two open questions of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n\nTo understand the connection between the world of factors of iids and descriptive combinatorics, we formalize the notion of toast construction used in both setups. We prove general results connecting the toast constructions with oblivious local algorithms and this enables us to understand the $\\mathsf{TAIL}$ complexity of a number of concrete problems like vertex $3$-coloring, edge $4$-coloring, perfect matchings and others. \n\nIn general, we hope that the concept of $\\mathsf{TAIL}$ complexities will find other usecases. We discuss in \\cref{subsec:congestLCA} how it relates to studied notions of oblivious algorithms, energy complexity, or its relevance to the $\\mathsf{LCA}$ model.\\todo{Oleg se ptal jestli uz ma tady vedet co je LCA a jestli to nekde neprehlid}\nIn fact, the $\\mathsf{TAIL}$ complexity approach was already considered e.g. in the celebrated MIS paper of \\cite{ghaffari2016MIS} (see the introduction of his paper) and the $\\mathsf{TAIL}$ bound on the MIS algorithm there may be the reason of its wide applicability. \n\nWe now list concrete results that we prove (or they follow from the work of others, see \\cref{fig:big_picture_grids}). \n\n\\paragraph{Classification Theorems for $\\mathsf{TAIL}$}\n\nWe prove several results that one can view as arrows in \\cref{fig:big_picture_grids}. Those can be viewed as an analogy to the line of work aiming for the classification of $\\mathsf{LOCAL}$ complexities discussed in \\cref{subsec:preliminaries_local}. \n\n\\begin{itemize}\n \\item $\\mathsf{TAIL}(\\mathrm{tower}(\\omega(r))) = \\mathsf{LOCAL}(O(1))$\\hspace*{1cm} We show that any oblivious local algorithm with tail decay faster than tower function can be in fact solved in $O(1)$ local complexity. This answers a question of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n \\item $\\mathsf{LOCAL}(O(\\log^* n)) = \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r))) = \\mathsf{CONTINUOUS}$\\hspace*{1cm} There are three natural definitions of classes of basic symmetry breaking problems in all three considered fields and they are the same \\cite{GJKS,HolroydSchrammWilson2017FinitaryColoring}. \n As a direct application one gets that from the three proofs that $3$-coloring of grids is not in $\\mathsf{LOCAL}(O(\\log^* n))$ \\cite{brandt_grids}, not in $ \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$ \\cite{HolroydSchrammWilson2017FinitaryColoring} and not in $ \\mathsf{CONTINUOUS}$ \\cite{GJKS} any one implies the others. These three classes are the same in a much more general setup (see \\cref{subsec:CONT}). \n \\item $\\mathsf{TAIL}(\\omega(r^{d-1})) = \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$\\hspace*{1cm} We prove that any problem that can be solved with an oblivious local algorithm with tail decay faster than $r^{d-1}$, can in fact be solved with a tower tail decay. This answers a question of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n \\item Similarly to the time hierarchy theorems of the $\\mathsf{LOCAL}$ model \\cite{chang2017time_hierarchy}, we study the range of possible $\\mathsf{TAIL}$ complexities and give examples of problems LCL problems with $\\mathsf{TAIL}$-complexities $\\Theta(r^{j - o(1)})$ for any $1 \\le j \\le d-1$. \n\\end{itemize}\n\n\\paragraph{$\\TOAST$ constructions}\nWe formally define $\\TOAST$ algorithms (we mentioned earlier that these constructions are used in descriptive combinatorics and the theory of random processes) and prove results about them. \nLet us first explain the toast constructions in a more detail. A $q$-toast in the infinite graph $\\mathbb{Z}^d$ is a collection $\\mathcal{D}$ of finite connected subsets of $\\mathbb{Z}^d$ that we call \\emph{pieces} such that (a) for any two nodes of the graph we can find a piece that contains both,\\todo{it is enough that every vertex is contained in a piece} (b) boundaries of any two different pieces are at least distance $q$ apart. \nThat is, imagine a collection of bigger and bigger pieces that form a laminar family (only inclusions, not intersections) that, in the limit, covers the whole infinite graph.\n\nThis structure can be built in the descriptive combinatorics setups \\cite{FirstToast,GJKS,Circle} but also with oblivious local algorithms \\cite{HolroydSchrammWilson2017FinitaryColoring,Spinka}. In fact, we show that one can construct the structure with $1\/r^{1-o(1)}$ $\\mathsf{TAIL}$ decay in \\cref{thm:toast_in_tail}. This is best possible up to the $o(1)$ loss in the exponent. \nIf one has an inductive procedure building the toast, we can try to solve an LCL problem during the construction. That is, whenever a new piece of the toast is constructed, we define how the solution of the problem looks like at this piece. This is what we call a $\\TOAST$ algorithm. \nThere is a randomized variant of the $\\TOAST$ algorithm and a lazy variant. In the lazy variant one can postpone the decision to color given piece, but it needs to be colored at some point in the future. \nNext, we prove several relations between the $\\TOAST$ classes (see \\cref{fig:big_picture_grids}). \n\n\\begin{itemize}\n \\item $\\TOAST \\subseteq \\mathsf{TAIL}(r^{1-o(1)})$\\hspace*{1cm} This result allows us to translate toast-algorithm constructions used in descriptive combinatorics to the $\\mathsf{TAIL}$-setting. This includes problems like vertex $3$-coloring (this answers a question of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}), edge $2d$-coloring, or perfect matching. \n \\item $\\RTOAST = \\TOAST$\\hspace*{1cm} This derandomization result for example implies on the first glance surprising inclusion $\\mathsf{FINDEP} \\subseteq \\mathsf{BOREL}$ -- these two classes are discussed in \\cref{sec:OtherClasses}. \n \\item $\\mathsf{TAIL} = \\RLTOAST$\\hspace*{1cm} Here, $\\mathsf{TAIL}$ is defined as the union of $\\mathsf{TAIL}(f)$ over all functions $f \\in \\omega(1)$. This result shows how the language of toast constructions can in fact characterize the whole hierarchy of constructions by oblivious local algorithms. \n \\item We show that the problem of $2d$ edge coloring of $\\mathbb{Z}^d$ is in $\\TOAST$, hence in $\\mathsf{BOREL}$ which answers the question of Gao et al \\cite{GJKS}. Independently, this was proven by \\cite{ArankaFeriLaci}. \n\\end{itemize}\n\n\\subsection{Roadmap}\n\nIn \\cref{sec:preliminaries} we recall the basic definitions regarding the $\\mathsf{LOCAL}$ model and explain the formal connection between finitary factors of iids. \nNext, in \\cref{sec:speedup_theorems} we prove the basic speedup theorems for $\\mathsf{TAIL}$. \nIn \\cref{sec:Toast} we define $\\TOAST$ algorithms and relate it to the $\\mathsf{TAIL}$ model. \nIn \\cref{sec:SpecificProblems} we discuss what our results mean for specific local problems and in \\cref{sec:Inter} we construct problems with $r^{j-o(1)}$ $\\mathsf{TAIL}$ complexities for $1 \\le j \\le d-1$. Finally, in \\cref{sec:OtherClasses} we discuss connection of our work to other classes of problems and finish with open questions in \\cref{sec:open_problems}. \n\n\n\n\n\n\\begin{comment}\nHonza - Introduction\n\nA \\emph{locally checkable problems (LCL) $\\Pi$ on a family of graphs $\\mathbf{G}$} is a procedure that checks whether a given ``candidate'' vertex coloring of $G\\in \\mathbf{G}$ satisfies a local constraints given by $\\Pi$\\footnote{\nFormally, $\\Pi$ is a triplet $\\Pi=(b,t,\\mathcal{P})$, where $b$ is a set of colors, $t\\in \\mathbb{N}$ is the locality of $\\Pi$ and $\\mathcal{P}$ is a set of allowed configurations around each vertex of locality at most $t$.}.\nMore concretely, given a graph $G=(V(G),E(G))\\in \\mathbf{G}$ and a ``candidate'' coloring $c:V(G)\\to b$, we say that $c$ is a \\emph{$\\Pi$-coloring} if for every $v\\in V(G)$ the restriction of $c$ to $t$-neighborhood of $v$ is in $\\mathcal{P}$, where $t$ is the locality of $\\Pi$ and $\\mathcal{P}$ is the set of constraints of $\\Pi$.\n\nIn this paper the family $\\mathbf{G}$ always refers to graphs that look locally like the standard Cayley graph of $\\mathbb{Z}^d$, $\\mathbb{Z}^d$, for some $d\\ge 1$.\nThis includes $\\mathbb{Z}^d$ as well as various finite $d$-dimensional tori with oriented edges that are labeled by the standard generators of $\\mathbb{Z}^d$.\nThe dimension $d$ is always assumed to be bigger than $1$, since for $d=1$ there is a full classification of LCL problems.\nWe omit mentioning ${\\bf G}$ and refer to all LCL problems on the family ${\\bf G}$ simply as \\emph{LCL problems (on $\n\\mathbb{Z}^d$)}.\nWhen we want to stress that a given LCL problem is on $\\mathbb{Z}^d$ we write $d$ in the superscript, i.e., $\\Pi^d$.\n\nWe study a complexity of LCL problems in various settings.\nIn the most abstract terms, a complexity of a given LCL problem $\\Pi$ is measured in terms of how difficult it is to \\emph{transform} a \\emph{given labeling} of a graph into a $\\Pi$-coloring.\nThe different contexts varies in what is considered to be an allowed starting labeling and what is allowed transformation.\nMost relevant for our investigation is the following probability notion.\nStart with a graph $G$ and assign to each vertex $v\\in V(G)$ a real number from $[0,1]$ uniformly and independently, i.e., \\emph{iid labeling} of $G$.\nA function $F$ that transforms these iid labels in a measurable and equivaraint way into a different coloring is called \\emph{factor of iid (fiid)}.\nHere equivariant refers to the group of all symmetries of $G$.\nIt is not hard to see that the group of all symmetries of $\\mathbb{Z}^d$ is the group $\\mathbb{Z}^d$ itself that acts by shift-translations.\nA fiid $F$ on $\\mathbb{Z}^d$ is called \\emph{finitary fiid (ffiid)} if for almost every iid labeling of $\\mathbb{Z}^d$ the color $F(x)({\\bf 0}_d)$ only depends on the labels of $x$ on some finite neighborhood of ${\\bf 0}_d$.\nThat is to say, if $x$ and $x'$ agrees on this neighborhood, then $F(x)({\\bf 0}_d)=F(x')({\\bf 0}_d)$.\nGiven ffiid $F$, we write $R_\\mathcal{A}$ for the random variable that codes minimum diameter of such neighborhood and we refer to it as the coding radius of $F$.\n\nWe say that a ffiid $F$ \\emph{solves an LCL problem $\\Pi$} if it produces a $\\Pi$-coloring for almost every iid labeling of $\\mathbb{Z}^d$.\nan LCL problem $\\Pi$ is in the class $\\mathsf{TAIL}$ if $\\Pi$ admits a ffiid solution.\nFollowing Holroyd, Schramm and Wilson \\cite{HolroydSchrammWilson2017FinitaryColoring}, the complexity of LCL problems in $\\mathsf{TAIL}$ can be measured on a finer scale by the properties of the coding radius $R_\\mathcal{A}$, where $F$ runs over all ffiid's that solve the given problem.\nSimilarly, we define when fiid $F$ solves $\\Pi$ and the class $\\mathsf{fiid}$.\n\nBefore we move to distributed algorithm we recall what is known about the notorious class of LCL problems that are the \\emph{proper vertex $k$-colorings of $\\mathbb{Z}^d$}, $\\Pi^d_{\\chi,k}$, where $k>1$, in the context of ffiid.\nIt is known that if $d>1$ and $k>3$, then $\\Pi^d_{\\chi,k}$ admits ffiid solution with a tower tail decay, see Holroyd et.al~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\nIf $d>0$ and $k=2$, then $\\Pi^d_{\\chi,2}$ is not even in $\\mathsf{fiid}$.\nThe most interesting case is when $d>1$ and $k=3$.\nIt was shown by Holroyd et al~\\cite{HolroydSchrammWilson2017FinitaryColoring} that in that case $\\Pi^d_{\\chi,3}$ is in the class $\\mathsf{TAIL}$ and every ffiid that solves $\\Pi^d_{\\chi,3}$ must have infinite the $2$-moment of its coding radius. \nMoreover, they found a ffiid that solve $\\Pi^d_{\\chi,3}$ and has finite some moments of its coding radius.\nHere we strengthen their result by showing that ${\\bf M}(\\Pi^d_{\\chi,3})=1$ and every ffiid solution of $\\Pi^d_{\\chi,3}$ has infinite the $1$-moment of its coding radius.\n\nIn the context of deterministic distributed computing the setting is the following.\nStart with finite graph $G$ on $n$-vertices (in our case $n$ is big enough and $G$ is a torus) and assign to every vertex of $G$ its unique identifier from $[n]$.\\footnote{In generality, the identifier could be assigned from some domain that has size polynomial in $n$.}\nA \\emph{local algorithm $\\mathcal{A}_n$ of locality $t_n$} is a rule that computes color of a vertex $v$ based only on its $t$-neighborhood and the unique identifiers in this neighborhood.\nWe say that $\\mathcal{A}_n$ \\emph{solves} $\\Pi$ if for every (admissible)\\footnote{As a torus $G$ needs to be wider than the locality of $\\mathcal{A}_n$} graph $G$ on $n$-vertices and any assignment of unique identifiers to $G$ the rule $\\mathcal{A}_n$ produces a $\\Pi$-coloring.\nWe say that $\\Pi$ is in the class $\\mathsf{LOCAL}(O(\\log^*(n)))$ if there is a sequence of local rules \\todo{co je to local rule?}$\\{\\mathcal{A}_n,t_n\\}_{n\\in \\mathbb{N}}$ such that $\\mathcal{A}_n$ has locality $t_n$, solves $\\Pi$ for every $n\\in\\mathbb{N}$ and $n\\mapsto t_n$ is in the class $O(\\log^*(n))$.\nThe reason why we picked $\\log^*(n)$ is that there is a complete classification of LCL problems on oriented labeled grids.\nRoughly speaking, an LCL problem $\\Pi$ is either \\emph{trivial}, i.e., there is a single color that produces $\\Pi$-coloring, \\emph{local}, i.e., there is an algorithm that needs locality $\\Theta(\\log^*(n))$, or \\emph{global}, i.e., every local algorithm that solves $\\Pi$ needs to see essentially the whole graph.\\footnote{Same classification holds in the context of randomized distributed algorithms.}\n\nWe use this classification to provide several speed-up results in the context of ffiid solutions.\nWe note that the connection between distributed computing and probability setting was previously observed by several authors.\nFor example, it follows easily from the result of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring} that local coloring problems that are not global in the distributed setting admit ffiid solution with tower tail decay.\nIn fact, Brandt et al. \\cite{brandt_grids} asked:\n\\begin{quote}\n\\emph{\nHowever, despite the fact that techniques seem to translate between finitary colourings\nand distributed complexity, it remains unclear how to directly translate results from one setting to\nthe other in a black-box manner; for instance, can we derive the lower bound for 3-colouring from\nthe results of Holroyd et al. [24], and does our complexity classification imply answers to the open\nquestions they pose?\n}\n\\end{quote}\nAs one of our main result we address their question.\n\nThe most abstract setting that we consider is Borel combinatorics.\nIn this context one studies graphs that additionally posses a completely metrizable and separable topological structure (or a standard Borel structure) on its vertex set.\nIt follows from the work of Seward and Tucker-Drob \\cite{ST} and Gao, Jackson, Krohne and Seward \\cite{GJKS} that the canonical object to investigate is the graph structure on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$, the free part of the shift action $\\mathbb{Z}^d\\curvearrowright 2^{\\mathbb{Z}^d}$, induced by the standard generating set of $\\mathbb{Z}^d$. \nThe topological and Borel structure on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ is the one inherited from the product topology on $2^{\\mathbb{Z}^d}$.\nWe say that an LCL problem $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ if there is a continuous and equivariant map from $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ into the space of $\\Pi$-colorings.\nThe class $\\mathsf{BOREL}$ is defined similarly, only the requirement on the map to be continuous is relaxed to be merely Borel measurable.\nBernshteyn \\cite{Bernshteyn2021LLL} proved that LCL problems that are not global in the distributed setting admit continuous and therefore Borel solution.\nIn another words, he showed that the class $\\mathsf{LOCAL}(O(\\log^*(n)))$ is contained in $\\mathsf{CONTINUOUS}$.\nIt is not hard to show that the other inclusion follows from the Twelve Tile Theorem of Gao et al \\cite{GJKS}, i.e., the class $\\mathsf{CONTINUOUS}$ in fact coincide with the class $\\mathsf{LOCAL}(O(\\log^*(n)))$.\nHere we investigate in greater detail what is the connection between existence of a Borel solution and a tail behavior of ffiid.\nIn fact all known Borel constructions uses the following structure that can be rephrased without mentioning topology or Borel measurability.\n\nA \\emph{$q$-toast on $\\mathbb{Z}^d$} is a collection $\\mathcal{D}$ of finite connected subsets of $\\mathbb{Z}^d$ such that\n\\begin{itemize}\n \\item for every $g,h\\in \\mathbb{Z}^d$ there is $D\\in\\mathcal{D}$ such that $g,h\\in D$,\n \\item if $D,E\\in \\mathcal{D}$ and $d(\\partial D,\\partial E)\\le q$, then $D=E$.\n\\end{itemize}\nWe say that an LCL problem $\\Pi$ is in the class $\\TOAST$ if there is $q\\in \\mathbb{N}$ and an equivariant rule $\\Omega$ that extends partial colorings such that if we apply $\\Omega$ on every $q$-toast $\\mathcal{D}$ inductively on the depth, then we produce a $\\Pi$-coloring.\nIt is not known whether the classes $\\mathsf{BOREL}$ and $\\TOAST$ coincide, however, it follows from the work of Gao et al.~\\cite{GJKS2} that $\\TOAST\\subseteq \\mathsf{BOREL}$.\nThis is proved by constructing a toast structure in a Borel way on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$, see Marks and Unger~\\cite{Circle} for application and proof of this fact.\nThe specific problems that belong to the class $\\TOAST$ include $3$-vertex coloring, perfect matching, $2d$-edge coloring, Schreier decoration (for non-oriented grids), some rectangular tilings, space filling curve etc.\nThese results were proved by many people in slightly different language, \\cite{GJKS2, ArankaFeriLaci, spencer_personal}.\nWe utilize here all the constructions and show that these all belong to $\\TOAST$.\n\nWe connect the classes $\\TOAST$ and $\\mathsf{TAIL}$ as follows.\nWe show that for every $q\\in \\mathbb{N}$ and $\\epsilon>0$ there is a ffiid $F$ that produces $q$-toast that has the $(1-\\epsilon)$-moment finite.\nA slightly stronger result implies that, in fact, ${\\bf M}(\\Pi)\\ge 1$, whenever an LCL problem $\\Pi$ is in the class $\\TOAST$.\nThis shows that ${\\bf M}(\\Pi)\\ge 1$ for all the problems mentioned above.\nWe note that in the special case when $d=2$ and the corresponding LCL problem $\\Pi$ does not admit ffiid solution with a tower tail decay, we deduce that ${\\bf M}(\\Pi)=1$.\nThis follows from our speed-up result in the context of ffiid.\n\nTraditionally there are two more classes that are investigated in Borel combinatorics, $\\mathsf{BAIRE}$ and $\\mathsf{MEASURE}$.\nTo define the class $\\mathsf{BAIRE}$ we replace the topological requirement on continuity by \\emph{Baire measurability}, i.e., we require a solution to be continuous only on a \\emph{comeager} set, see the work of Bernshteyn \\cite{Bernshteyn2021LLL}.\nThe class $\\mathsf{MEASURE}$ generalizes the setting of fiid.\nThe goal here is to find a Borel solution on a conull set for a given Borel graph that is endowed with a Borel probability measure.\nThis direction has been recently investigated by Bernshteyn \\cite{Bernshteyn2021LLL} for general graphs, where he proved a measurable version of Lovasz's Local Lemma.\nSee \\cref{fig:big_picture_grids} for all the relations that are known to us.\n\\begin{figure}\n \\centering\n \\includegraphics[width=\\textwidth]{fig\/big_picture_grids.png}\n \\caption{Big picture of lcl for $\\mathbb{Z}^d$}\n \n\\end{figure}\n\n\n\n\n\n\n\n\\subsection{The Results}\n\nIn this subsection we formulate our main results.\nThe formulation uses mostly the language of ffiid because the tail behavior of the coding radius allows to compare problems on a quantitative and, as it turns out, rich scale.\nOne should always re-interpret the result with the help of Figure~\\ref{fig:big_picture_grids} to different contexts.\nIronically enough, the techniques that are most valuable when investigating ffiid solutions come either from local algorithms or Borel combinatorics, where such a rich and quantitative distinction is not possible.\nRecall that in the distributed algorithms there is a trichotomy of LCL problems, i.e., trivial, local or global, and in Borel combinatorics we distinguish mainly between trivial, continuous, Borel, and global problems.\n\n\nWe start with results that use ideas from distributed algorithms.\nFirst, we prove that LCL problems that are not global in distributed algorithms are precisely problems that admit ffiid solution with a tower tail decay.\n\n\\begin{theorem*}\nLet $d\\in \\mathbb{N}$ and $\\Pi$ be an LCL problem on $\\mathbb{Z}^d$.\nThen $\\Pi$ is not global in the local world, i.e., there is a deterministic algorithm in the local model that runs in $O(\\log^*(n))$ rounds, if and only if there is an ffiid solution $F$ of $\\Pi$ such that\n$$\\P(R_\\mathcal{A}>r)\\le \\frac{1}{tower(Cr)}$$\nfor some $C>0$ and large enough $r\\in \\mathbb{N}$.\n\\end{theorem*}\n\nThe translation between ffiid and distributed algorithms allows to deduce the following result that answers problem (ii) of Holroyd at el.~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\n\n\\begin{theorem*}\nLet $d\\in \\mathbb{N}$ and $\\Pi$ be an LCL problem on $\\mathbb{Z}^d$.\nIf $\\Pi$ admits an ffiid solution $F$ that has faster tail decay than any tower function, then $\\Pi$ is trivial, i.e., there is a color such that a constant coloring with this color is a correct solution.\n\\end{theorem*}\n\nAs a last application of ideas from the distributed algorithms, we show the following speed-up theorem for ffiid.\n\n\\begin{theorem*}\nLet $d,t\\in \\mathbb{N}$.\nThere is $C_{d,t}>0$ such that if an LCL problem $\\Pi$ admits an ffiid solution $F$ such that\n$$\\P(R_\\mathcal{A}>r)\\le C_{d,t}r^{d-1},$$\nthen there is a deterministic algorithm in the local model that runs in $O(\\log^*(n))$ rounds and solves $\\Pi$.\n\\end{theorem*}\n\nTo understand better LCL problems that are global in the distributed algorithms or admit no continuous solution in the context of Borel combinatorics we introduce the following numerical invariant.\nThe inspiration comes from the paper of Holroyd et al.~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\n\n\\begin{definition*}\nLet $\\Pi$ be an LCL on $\\mathbb{Z}^d$.\nWe define ${\\bf M}(\\Pi)\\in [0,+\\infty]$ to be the supremum over all $\\alpha\\in [0,+\\infty)$ such that there is an ffiid solution $F$ of $\\Pi$ such that $R_F$ has finite the $\\alpha$-moment.\n\\end{definition*}\n\nCombination of previous results immediately implies the following strengthening of the upper bound on the necessary infinite moment of ffiid solution of the proper vertex $3$-coloring problem, $\\Pi^d_{\\chi,3}$, from Holroyd et al.~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\nMoreover, this answers their problem (iii).\n\n\\begin{theorem*}\nLet $d\\in \\mathbb{N}$ and $\\Pi$ be an LCL problem on $\\mathbb{Z}^d$.\nIf ${\\bf M}(\\Pi)>d-1$, then $\\Pi$ admits an ffiid solution with tower tail decay.\n\\end{theorem*}\n\nAs a first application of ideas from Borel combinatorics, i.e., the toast structure, we compute the optimal tail decay for $\\Pi^d_{\\chi,3}$, thus answering a problem (i) of Holroyd et al.~\\cite{HolroydSchrammWilson2017FinitaryColoring}.\n\n\\begin{theorem*}\nLet $d>1$ and $\\Pi^d_{\\chi,3}$ be the proper vertex $3$-coloring problem.\nThen ${\\bf M}(\\Pi^d_{\\chi,3})=1$.\n\\end{theorem*}\n\nA high level explanation of a toast construction in the context of ffiid says that if $r\\in \\mathbb{N}$ is given, then the fraction of vertices not covered by first $r$-layers of the toast is at most $\\frac{1}{r}$.\nSimilar construction that uses only $1$-dimensional skeleton of some rectangular tiling of $\\mathbb{Z}^d$ improve this upper bound to $\\frac{1}{r^{d-1}}$.\nWe show that this type of structure is enough to find a perfect matching in a bigger graph then $\\mathbb{Z}^d$, a graph where we moreover add diagonal edges in every $2$-dimensional hyperplane.\nThis yields an existence of intermediate problems from the point of view of ${\\bf M}$.\n\n\\begin{theorem*}\nLet $2\\le \\ell\\le d\\in \\mathbb{N}$.\nThen there is an LCL problem $\\Pi$ on $\\mathbb{Z}^d$ that satisfies ${\\bf M}(\\Pi)=\\ell-1$.\n\\end{theorem*}\n\nAs a last result we summarize what is known about specific problems.\nSince we work mainly with oriented grids we state everything in this context, however, some of the results can be deduced for non-oriented grids as well, see Section~\\ref{sec:SpecificProblems}.\nIn particular, we answer Question~6.5 of Gao et al.~\\cite{GJKS} about Borel edge chroamtic number of $\\operatorname{Free}(2^{\\mathbb{Z}^d})$.\n\n\\begin{theorem*}\nLet $d\\ge 2$ and $\\Pi$ be any of the following LCL problems:\n\\begin{enumerate}\n \\item proper vertex $3$-coloring, \\cite{GJKS}, \\cite{spencer_personal},\n \\item proper edge $2d$-coloring,\n \\item rectangular tiling with rectangles from some set $\\mathcal{R}$ that satisfies \\cref{eq:GCD}, \\cite{Unger},\n \\item perfect matching, \\cite{GJKS}.\n\\end{enumerate}\nThen $\\Pi$ is in the class $\\TOAST$.\nConsequently, $\\Pi$ is in the class $\\mathsf{BOREL}$ and satisfies $1\\le {\\bf M}(\\Pi)$.\n\\end{theorem*}\n\n\n\nThe paper is organized as follows...\n\n\n\\end{comment}\n\n\n\n\n\n\n\\section{Preliminaries}\n\\label{sec:preliminaries}\n\nLet $d\\in \\mathbb{N}$ be fixed.\nThe class of graphs that we consider in this work contains mostly the infinite grid $\\mathbb{Z}^d$, that is $\\mathbb{Z}^d$ with the Cayley graph structure given by the standard set of generators \n$$S=\\{(1,0,\\dots, 0),(0,1,\\dots, 0), \\dots (0,\\dots,0,1)\\}. $$\nWe also used the finite variant of the graph: the torus of side length $\\Theta(n^{1\/d})$ with edges oriented and labeled by $S$.\nUsually, when we talk about a graph $G$ we mean one of the graphs above and it follows from the context whether we mean finite or infinite graph. \n\nGiven a graph $G$, we denote as $V(G)$ the set of vertices, as $E(G)$ the set of edges and as $d_G$ the \\emph{graph distance} on $G$, i.e., $d_G(v,w)$ is the minimal number of edges that connects $v$ and $w$.\nWe write $\\mathcal{B}_G(v,r)$ for the \\emph{ball of radius $r\\in \\mathbb{N}$} around $v\\in V(G)$. When $G$ is the infinite grid $\\mathbb{Z}^d$, then we leave the subscripts blank.\nA \\emph{rooted graph} is a pair $(G,v)$ where $G$ is a graph and $v\\in V(G)$.\nWe denote as $[G,v]$ the \\emph{isomorphism type} of the rooted graph $(G,v)$, where $(G,v)$ and $(G',v')$ are isomorphic if there is a bijection $\\varphi:V(G)\\to V(G')$ that sends $v$ to $v'$ and preserves edge orientation and labeling.\n\nLet $D$ be a set, $G$ a graph and $c:V(G)\\to D$ be a labeling of vertices.\nThen we write $G_c=(G,c)$ for the pair graph + labeling.\nSimilarly we denote as $(G_c,v)$ the rooted graph $(G,v)$ with the labeling $c$, $[G_c,v]$ its isomorphism class and $\\mathcal{B}_{G_c}(v,r)$ the labeled ball of radius $r\\in \\mathbb{N}$ around $v$, where $v\\in V(G)$.\nOften we do not mention the labels and tacitly assume that they are understood from the context.\nAlso, the label set might depend on the size of of the vertex set, e.g., unique labels from $[|V(G)|]=\\{1,\\dots,|V(G)|\\}$ etc.\n\n\\begin{comment}\n\nWe always consider the group $\\mathbb{Z}^d$ with the standard set of generators $S$, i.e.,\n$$S=\\{(1,0,\\dots, 0),(0,1,\\dots, 0), \\dots (0,\\dots,0,1)\\}.$$\nThe \\emph{Cayley graph} $\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right):=\\operatorname{Cay}\\left(\\mathbb{Z}^d,S\\right)$ of $\\mathbb{Z}^d$ is then defined on the vertex set $\\mathbb{Z}^d$ with oriented and $S$-labeled edges \n$$\\left\\{(g,s+g):s\\in S, g\\in \\mathbb{Z}^d\\right\\}=E\\left(\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)\\right).$$\nWe denote the identity in $\\mathbb{Z}^d$ as ${\\bf 0}_d$. \nWhen we talk about $\\mathbb{Z}^d$ next, we mean formally the appropriate Cayley graph with labelled edges, that is, we think of the graph together with its consistent orientation. \n\n\nA \\emph{graph} $G$ is a pair $(V(G),E(G))$, where $V(G)$ is a vertex set and $E(G)$ consists of oriented and $S$-labeled edges.\nGiven a graph $G$, we denote as $d_G$ the \\emph{graph distance} on $G$, i.e., $d_G(v,w)$ is the minimal number of edges that connects $v$ and $w$.\nWe write $\\mathcal{B}_G(v,r)$ for the \\emph{ball of radius $r\\in \\mathbb{N}$} around $v\\in V(G)$. Usually, $G$ is just the grid $\\mathbb{Z}^d$, then we leave the subscript blank. \nA \\emph{rooted graph} is a pair $(G,v)$ where $G$ is a graph and $v\\in V(G)$.\nWe denote as $[G,v]$ the isomorphism type of the rooted graph $(G,v)$, where $(G,v)$ and $(G',v')$ are isomorphic if there is a bijection $\\varphi:V(G)\\to V(G')$ that sends $v$ to $v'$ and preserves edge orientation and labeling.\n\n\\todo{vsechny strany toru maji delku n na 1\/d; zjevne vsechny problemy maji lokalitu nejvyse n na 1\/d}\n\n\\end{comment}\n\n\n\\begin{comment}\nWe write $\\operatorname{\\bf FG}^{S}$ for the set of all graphs that look in their $2$-neighborhoods as $\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)$, i.e., a graph $G=(V(G),E(G))$ is in $\\operatorname{\\bf FG}^{S}$ if $[\\mathcal{B}_G(v,2)]=\\left[\\mathcal{B}_{\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)}({\\bf 0}_d,2)\\right]$ for every $v\\in V(G)$.\nNote that $\\mathbb{Z}^d\\in \\operatorname{\\bf FG}^{S}$ and there are finite graphs in $\\operatorname{\\bf FG}^{S}$, as an example take any torus that is big enough.\nWrite $\\operatorname{\\bf FG}_{\\bullet}^{S}$ for the rooted graphs of the form\n$$(\\mathcal{B}_G(v,r),v)$$\nwhere $G\\in \\operatorname{\\bf FG}^{S}$, $v\\in V(G)$ and $r\\in \\mathbb{N}$.\nWe say that a rooted graph $(G,v)\\in \\operatorname{\\bf FG}_{\\bullet}^{S}$ has \\emph{diameter at most $m\\in \\mathbb{N}$} if $(\\mathcal{B}_G(v,m),v)=(G,v)$.\nNote that, in particular, we have $\\mathcal{B}_{\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)}({\\bf 0}_d,r)\\in \\operatorname{\\bf FG}_{\\bullet}^{S}$ for every $r\\in \\mathbb{N}$.\n\nLet $D$ be a set, $G$ a graph and $c:V(G)\\to D$ be a labeling of vertices.\nThen we write $G_c=(G,c)$ for the pair graph + labeling.\nSimilarly we denote as $(G_c,v)$ the rooted graph $(G,v)$ with the labeling $c$, $[G_c,v]$ its isomorphism class and $\\mathcal{B}_{G_c}(v,r)$ the labeled ball of radius $r\\in \\mathbb{N}$ around $v$, where $v\\in V(G)$.\nOften the label set $D$ depends on the size of of the vertex set, e.g., unique labels from $[|V(G)|]=\\{1,\\dots,|V(G)|\\}$ etc.\n\\end{comment}\n\n\n\n\\subsection{LCLs on $\\mathbb{Z}^d$}\n\nOur goal is to understand locally checkable problems (LCL) on $\\mathbb{Z}^d$. \n\n\\begin{definition}[Locally checkable problems \\cite{naorstockmeyer}]\nA \\emph{locally checkable problems (LCL)} on $\\mathbb{Z}^d$ is a quadruple $\\Pi=(\\Sigma_{in}, \\Sigma_{out},t,\\mathcal{P})$ where $\\Sigma_{in}$ and $\\Sigma_{out}$ are finite sets, $t\\in \\mathbb{N}$ and $\\mathcal{P}$ is an assignment\n$$[H_c,w]\\mapsto \\mathcal{P}([H_c,w])\\in \\{0,1\\},$$\nwhere $(H,w)$ is a rooted grid of radius at most $t$ together with an input and output coloring $c: V(H)\\to \\Sigma_{in}\\times\\Sigma_{out}$. \n \nGiven a graph $G$ together with a coloring $c:V(G)\\to \\Sigma_{in}\\times\\Sigma_{out}$, we say that $c$ is a \\emph{$\\Pi$-coloring} if\n$$\\mathcal{P}([\\mathcal{B}_{G_c}(v,t),v])=1$$\nfor every $v\\in V(G)$.\n\\end{definition}\n\n\n\n\\begin{comment}\n\nClasses of functions\n\nBefore we proceed to definition of local algorithms and their complexity we recall here the relevant notions concerning functions from $\\mathbb{N}$ to $\\mathbb{N}$. \nImportantly, we often use asymptotic notation as usual in Computer Science, e.g., we define the class of functions $O(f(n))$ as the union over all functions $g:\\mathbb{N} \\rightarrow \\mathbb{N}$ such that there exists $C > 0$ such that for all $n$ we have $g(n) \\le C f(n)$. Other symbols $\\Omega, \\Theta, o, \\omega$ are defined analogously and their usage is analogous (see e.g. \\cite{?}). \n\\end{comment}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{LOCAL model}\n\\label{subsec:preliminaries_local}\n\nHere we formally introduce the $\\mathsf{LOCAL}$ model for both deterministic and randomized algorithms and state the classification of LCLs on $\\mathbb{Z}^d$ from the $\\mathsf{LOCAL}$ point of view.\nInformally, a local algorithm $\\mathcal{A}$ of radius $r\\in \\mathbb{N}$ is a rule that takes a labeled graph $G_c$ and then re-labels each vertex $v\\in V(G)$ depending only on the isomorphism type of its neighborhood $[\\mathcal{B}_{G_c}(v,r),v]$ with a new label from some finite set $b$.\n\n\\paragraph{Local algorithms}\n\nHere, we formally define the complexity classes $\\mathsf{LOCAL}(\\mathcal{F})$ of deterministic local algorithms, where $\\mathcal{F}$ is a class of functions from $\\mathbb{N}$ to $\\mathbb{N}$. \nThe somewhat unintuitive part of the definition is that the algorithm takes as input \\emph{unique identifiers}. Usually, they are taken from some range of values that is polynomial in $n$, the number of vertices of the input graph. However, in most scenarios this is equivalent to the assumption that they are from range $\\{1,\\dots n\\}$, i.e., every identifier is used exactly once. We also use this definition here, for simplicity.\nIn all forthcoming definition we have some fixed LCL $\\Pi$ in mind.\n\n\\begin{definition}[A deterministic local algorithm for $\\mathbb{Z}^d$]\n\\label{def:local_algorithm\nLet $\\Sigma_{out}$ be a finite set.\nA \\emph{local (deterministic) algorithm $\\mathcal{A}$} is a sequence of pairs $\\{\\mathcal{A}_n,\\ell_n\\}_{n\\in\\mathbb{N}}$, where $\\ell_n\\in \\mathbb{N}$ and $\\mathcal{A}_n$ is an assignment, such that\n$$[\\mathcal{B}_{G_\\ell}(v,\\ell_n),v]\\mapsto \\mathcal{A}_n([\\mathcal{B}_{G_\\ell}(v,\\ell_n),v])\\in \\Sigma_{out}$$\nfor every $\\Sigma_{in}$-labeled graph $G$ (torus) on $n$ vertices, $v\\in V(G)$ and a bijection $\\ell:V(G)\\to [n]$.\nIn another words, $\\mathcal{A}_n$ is a function that is defined on isomorphism classes of rooted graphs of diameter at most $\\ell_n$ with injectively labeled vertices from the set $\\{1,\\dots, n\\}$ and with range $b$. \n\nGiven such a $G$ and a bijection $\\ell:V(G)\\to |V(G)|$, we define\n$$\\mathcal{A}(G_\\ell)(v):=\\mathcal{A}_{|V(G)|}([\\mathcal{B}_{G_\\ell}(v,\\ell_n),v])$$\nfor every $v\\in V(G)$.\n\\end{definition}\n\nNext we define the deterministic complexity of a given local coloring problem $\\Pi$.\n\n\\begin{definition}[Complexity of deterministic local algorithm for $\\mathbb{Z}^d$]\nWe say that a local algorithm $\\mathcal{A}=\\{\\mathcal{A}_n,\\ell_n\\}_{n\\in \\mathbb{N}}$ \\emph{solves} an LCL $\\Pi=(\\Sigma_{in},\\Sigma_{out},t,\\mathcal{P})$ if $\\mathcal{A}(G_l)$ is a $\\Pi$-coloring of $G$ for every finite $\\Sigma_{in}$-labeled graph $G$ and any bijection $\\ell:V(G)\\to |V(G)|$.\n\nWe say that $\\Pi$ is in the class $\\mathsf{LOCAL}(\\mathcal{F})$, where $\\mathcal{F}$ is a class of functions, if $n\\mapsto \\ell_n$ is in $\\mathcal{F}$.\n\\end{definition}\n\nA randomized local algorithm and its complexity is defined analogously but instead of labeling with unique identifiers each node has an iid random string and the algorithm succeeds with probability $1 - 1\/n$.\n\nCrucial for our investigation are so-called \\emph{oblivious local algorithms}.\nRoughly speaking an algorithm is \\emph{oblivious} if it does not have knowledge of $n$ (see \\cite{Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring} for details). \nSince we do not consider deterministic oblivious algorithms, by an oblivious local algorithm we mean an oblivious randomized local algorithm.\n\n\\begin{definition}[Oblivious local algorithm]\n\\label{def:uniform_local}\nA \\emph{oblivious local algorithm} $\\mathcal{A}$ is a pair $(\\mathcal{C},f)$, where $\\mathcal{C}$ is a (measurable) collection of rooted graphs labeled with $\\Sigma_{in}$ and real numbers from $[0,1]$, and $f: \\mathcal{C} \\rightarrow \\Sigma_{out}$ is a measurable assignment.\n\nWe define $\\mathcal{A}(G)$ to be the labeling computed by $\\mathcal{A}$ on $G$, where for a given node $u$ we have $\\mathcal{A}(G)(u)=g([B_G(u, t), u])$ where $t$ is the smallest integer such that $[B_G(u, t),u] \\in \\mathcal{C}$.\n\\end{definition}\n\nNote that since there is no $n$ in the definition of oblivious local algorithm, the algorithm is correct if it always finishes with probability $1$, while randomized algorithms can have small probability of error dependent on $n$. Hence, one should think about such algorithms as an analogue of Las Vegas algorithms, while the classical randomized local algorithms correspond to Monte Carlo algorithms. \nAn advantage of not knowing $n$ is that it is possible to run $\\mathcal{A}$ on, e.g., the infinite grid $\\mathbb{Z}^d$ instead of just a torus. Usually, these two models give the same results but for toast algorithms that we consider later, there is a difference: While we define the algorithms for an infinite grid, for example the $3$-coloring algorithm makes sense also for a finite torus, while the perfect matching algorithm does not make sense in the finite world for the simple reason that the input graph can have an odd number of vertices. \n\n\n\n\n\n\n\n\n\\begin{comment}\n\\paragraph{Randomized local algorithms}\\todo{merge with previous paragraph}\n\nThe main difference between deterministic and randomized model is that the input of the local rule in the randomized case are graphs with vertices labeled uniformly independently at random by elements of the interval $[0,1]$ and we require that the output coloring is correct everywhere with high probability.\nThe complexity of local algorithm is then defined as the locality it needs.\n\n\\begin{definition}[Randomized local algorithm for $\\mathbb{Z}^d$]\nLet $b$ be a set.\nWe say that $\\mathbb{A}=\\{\\mathbb{A}_n,\\ell_n\\}_{n\\in\\mathbb{N}}$ is a randomized local algorithm for the class $\\operatorname{\\bf FG}^{S}$ if $\\ell_n\\in \\mathbb{N}$ and \n$$[\\mathcal{B}_{G_l}(v,\\ell_n),v]\\mapsto \\mathbb{A}_n([\\mathcal{B}_{G_l}(v,\\ell_n),v])\\in b$$\nis a measurable assignment, where $G\\in \\operatorname{\\bf FG}^{S}$, $|V(G)|=n$, $v\\in V(G)$ and $\\ell:V(G)\\to [0,1]$.\nIn another words, $\\mathbb{A}_n$ is a measurable function that is defined on isomorphism classes of rooted graphs of diameter at most $\\ell_n$ with vertices labeled by $[0,1]$ and with range $b$. \n\nGiven $G\\in \\operatorname{\\bf FG}^{S}$ such that $n=|V(G)|$ and a map $\\ell:V(G)\\to [0,1]$, we define\n$$\\mathbb{A}(G_l)(v):=\\mathbb{A}_n([\\mathcal{B}_{G_l}(v,\\ell_n),v])$$\nfor every $v\\in V(G)$.\n\\end{definition}\n\nNext we define when a randomized algorithm solves an LCL $\\Pi$ and the randomized complexity classes.\n\n\\begin{definition}[Complexity of randomized local algorithm for $\\mathbb{Z}^d$]\nLet $\\Pi=(b,t,\\mathcal{P})$ be an LCL.\nWe say that a randomized algorithm $\\mathbb{A}=\\{\\mathbb{A}_n,\\ell_n\\}_{n\\in \\mathbb{N}}$ solves $\\Pi$ if \n$$\\lambda^{\\oplus n}\\left(\\left\\{\\ell\\in [0,1]^{V(G)}:\\mathbb{A}(G_l) \\ \\operatorname{is \\ a} \\ \\Pi\\operatorname{-coloring} \\right\\}\\right)>1-\\frac{1}{n}$$\nholds for every $n\\in \\mathbb{N}$ and $G\\in \\operatorname{\\bf FG}^{S}$ such that $|V(G)|=n$ and\n$$[\\mathcal{B}_{G}(v,\\ell_n+t),v]=\\left[\\mathcal{B}({\\bf 0}_d,\\ell_n+t),{\\bf 0}_d\\right]$$\nfor every $v\\in V(G)$.\n\nWe say that $\\Pi$ is in the class $\\mathsf{R-LOCAL}(\\mathcal{F})$, where $\\mathcal{F}$ is a class of functions, if $n\\mapsto \\ell_n$ is in $\\mathcal{F}$.\n\\end{definition}\n\\end{comment}\n\n\\begin{comment}\n\\label{sec:clocal_definition}\n\nThe last class of local algorithm that we introduce serves on one hand as a main technical tool for deterministic speed up tricks and on the other hand allows a smooth translation of local algorithm into the measurable setting.\nIntuitively, an algorithm is in the class $\\mathsf{C-LOCAL}$ if it only needs a fixed finite locality (not depending on the size of a graph) given some coloring of the vertices with fixed number of colors (not depending on the size of a graph).\n\nLet $s\\in \\mathbb{N}$ and $G$ be a graph.\nDenote as $G^s$ the $s$-power graph of $G$, i.e., $V(G^s)=V(G)$ and $(x,y)\\in E(G^s)$ if and only if $0 r),$$\nwhat we mean is the following. \nFirst, for a given rooted neighborhood $(G,v)$ with a fixed $\\Sigma_{in}$-labeling we ask what is the probability that, after labeling $G$ with real numbers uniformly at random, that $(G,v)$ is in $\\mathcal{C}$. \nSecond, we take the supremum over all such values over all $\\Sigma_{in}$-labels.\n\n\\begin{definition}[$\\mathsf{TAIL}(f)$]\n\\label{def:tail}\nWe say that an LCL $\\Pi$ is in the class $\\mathsf{TAIL}(f)$ for some function $f:\\mathbb{N}\\to \\mathbb{N}$ if there is an oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi$ such that\n$$\\P(R_\\mathcal{A} > r)\\le \\frac{1}{f(r)}$$\nfor every large enough $r\\in \\mathbb{N}$. \n\nWe say that $\\Pi \\in \\mathsf{TAIL}$ if there exists $f \\in \\omega(1)$ such that $\\Pi \\in \\mathsf{TAIL}(f)$. \nWe say that $\\Pi$ is solved by a \\emph{block factor} if there is an oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi$ such that for some $r_0$ we have $\\P(R_\\mathcal{A} > r_0) = 0$.\n\\end{definition}\n\nHere things are made simple because we think of the algorithm as being run on an infinite graph so nothing depends on $n$ in this complexity measure (while it is a crucial parameter in the classical local complexity definition). \n\nWe also use the asymptotic notation when talking about $\\mathsf{TAIL}$ complexities. For example, an important class $\\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$ is defined as the class of LCLs $\\Pi$ such that there is a constant $C\\in \\mathbb{N}$ and oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi$ that satisfies\n$$\\P(R_\\mathcal{A}> r)\\le \\frac{1}{\\mathrm{tower}(Cr)}$$\nwhere $\\mathrm{tower}(x)=\\exp\\circ\\dots\\circ \\exp(1)$ is $\\lfloor x\\rfloor$-many times iterated exponential.\n\nA coarser way to measure the complexity of a problem $\\Pi$ is to ask what moments of $R_\\mathcal{A}$ could be finite, when $\\mathcal{A}$ is an oblivious local algorithm.\nFor example, the first moment finite corresponds to the fact that average time until an algorithm finishes is defined (for finite graphs this is independent on the size of the graph). \nIn particular, we investigate the question what is the least infinite moment for any oblivious local algorithm to solve a given problem $\\Pi$, as this language was used in \\cite{HolroydSchrammWilson2017FinitaryColoring}. \nRecall that $R_\\mathcal{A}$ has a \\emph{finite $\\alpha$-moment}, where $\\alpha\\in [0,+\\infty)$, if\n$$\\mathbb{E}\\left(\\left(R_\\mathcal{A}\\right)^\\alpha\\right)=\\sum_{r \\in \\mathbb{N}} r^\\alpha \\cdot \\P(R_\\mathcal{A} = r) <+\\infty.$$\n\n\\begin{definition}\nLet $\\Pi$ be an LCL.\nWe define ${\\bf M}(\\Pi)\\in [0,+\\infty]$ to be the supremum over all $\\alpha\\in [0,+\\infty)$ such that there is an oblivious $\\mathcal{A}$ that solves $\\Pi$ and has the $\\alpha$-moment of $R_\\mathcal{A}$ finite.\n\\end{definition}\n\nThe $\\alpha$ moment being finite roughly corresponds to $\\Omega(r^{\\alpha})$ $\\mathsf{TAIL}$ complexity. In fact we now recall the following observation.\n\n\\begin{fact}\n\\label{fact:moment_implies_tail}\nSuppose that the coding radius of an oblivious local algorithm $\\mathcal{A}$ has the $\\alpha$-moment finite.\nThen for every $C>0$ we have\n$$\\P(R_\\mathcal{A} > r) \\le C\\frac{1}{r^\\alpha}$$\nfor large enough $r\\in \\mathbb{N}$.\n\\end{fact}\n\\begin{proof}\nWe prove the contrapositive. Suppose the tail of the coding radius $R$ satisfies $\\P(R > r_i) > C \\cdot \\frac{1}{r^\\alpha_i}$ for some $C > 0$ and an infinite sequence $(r_i)_{i\\in \\mathbb{N}}$. \nAs the sequence $\\P(R> r_i)$ goes to zero, we may pass to a subsequence such that for each $i\\in \\mathbb{N}$ we have $\\P(R > r_{i+1}) \\le \\frac{1}{2}\\P(R > r_i)$. We call the new subsequence still $(r_i)_{i\\in \\mathbb{N}}$. \nWe have \n\\begin{align*}\n\\mathbb{E}(R^\\alpha) &\\ge \\sum_{i \\in \\mathbb{N}} r_i^\\alpha \\cdot (\\P(R> r_i) - \\P(R > r_{i+1}))\\\\\n&\\ge \\sum_{i \\in \\mathbb{N}} r_i^\\alpha \\cdot \n\\frac{1}{2}\\P(R > r_i)\n\\ge \\sum_{i \\in \\mathbb{N}} r_i^\\alpha \\cdot \n\\frac{C}{2} \\cdot \\frac{1}{r_i^\\alpha}\n= +\\infty\n\\end{align*}\nwhich finishes the proof. \n\\end{proof}\n\nFinally, an important fact is that a small $\\mathsf{TAIL}$ complexity of an oblivious randomized local algorithm implies that it has small randomized complexity (on finite graphs). This is because if we run an oblivious algorithm with tail decay $1\/f(r)$ on an $n$-node graph and mark its run as a failure if the coding radius exceeds $f^{-1}(n^2)$, the failure at some node happens only with probability $1\/n$. This was formally proven in \\cite{brandt_grebik_grunau_rozhon2021classification_of_lcls_trees_and_grids} and here we only state the result. \n\n\\begin{lemma}\n\\label{lem:tail->local}\nLet $\\mathcal{A}$ be an oblivious local randomized algorithm with $\\mathsf{TAIL}$ complexity $f(r)$. Then the randomized complexity of $\\mathcal{A}$ is $O(f^{-1}(n^2))$. \n\\end{lemma}\n\nIn fact, one can prove that for $o(\\log n)$ local complexities the two notions are the same. We refer to the reader to the discussion in \\cite{brandt_grebik_grunau_rozhon2021classification_of_lcls_trees_and_grids} for more details. \n\n\\begin{theorem}[\\cite{brandt_grebik_grunau_rozhon2021classification_of_lcls_trees_and_grids}]\n\\label{thm:tail=uniform}\nLet $\\mathcal{A}$ be an oblivious local randomized algorithm with randomized complexity $f(n) = o(\\log n)$. Then this algorithm in fact never returns an incorrect solution and its tail complexity is $\\Omega(f^{-1}(r))$. \n\\end{theorem}\n\n\n\n\n\n\n\n\n\n\\subsection{Subshifts of finite type}\n\\label{subsec:subshifts}\n\nHere we connect our definition of LCLs and oblivious local algorithms with the well-studied notions of subshifts of finite type and finitary factors of iid labelings.\n\nTo define subshifts of finite type we restrict ourselves to LCLs without input.\nThat is $\\Sigma_{in}$ is empty or contains only one element.\nMoreover, for convenience we set $b:=\\Sigma_{out}$.\nThen an LCL $\\Pi$ is of the form $\\Pi=(b,t,\\mathcal{P})$.\nWe define $X_\\Pi\\subseteq b^{\\mathbb{Z}^d}$ to be the space of all $\\Pi$-colorings.\nThat is we consider all possible labelings of $\\mathbb{Z}^d$ by elements from $b$ and then include in $X_{\\Pi}$ only those that satisfies the finite constraints $\\mathcal{P}$.\nViewing $\\mathbb{Z}^d$ as a group, there is a natural shift action of $\\mathbb{Z}^d$ on (the Cayley graph) $\\mathbb{Z}^d$.\nNamely, $g\\in \\mathbb{Z}^d$ shifts $h$ to $g+h$.\nThis action naturally ascend to an action $\\cdot $ on $b^{\\mathbb{Z}^d}$ as\n$$(g\\cdot x)(h)=x(g+h)$$\nwhere $x\\in b^{\\mathbb{Z}^d}$.\nIt is not hard to see (since $\\mathcal{P}$ is defined on isomorphism types and not rooted in any particular vertex) that $X_\\Pi$ is invariant under the shift map.\nThen we say that $X_\\Pi$ together with the shift structure is a \\emph{subshift of finite type}.\n\nA \\emph{finitary factor of iid process (ffiid) $F$} is a function that takes as an input $\\mathbb{Z}^d$ labeled with real numbers from $[0,1]$, independently according to Lebesgue measure, and outputs a different labeling of $\\mathbb{Z}^d$, say with elements of some set $b$. \nThis labeling satisfies the following: (a) $F$ is equivariant, that is for every $g\\in \\mathbb{Z}^d$ shifting the real labels with $g$ and applying $F$ is the same as applying $F$ and then shifting with $g$, (b) the value of $F(x)({\\bf 0})$ depends almost surely on the real labels on some random but finite neighborhood of ${\\bf 0}$.\nSimilarly as for oblivious local algorithms one defines the random variable $R_F$ that encodes minimal such neighborhood.\nWe say that ffiid $F$ solves $\\Pi$ if $F(x)\\in X_{\\Pi}$ with probability $1$.\n\nThe difference between ffiid and oblivious local algorithms is just formal.\nNamely, given an oblivious local algorithm $\\mathcal{A}$, we can turn it into an ffiid $F_\\mathcal{A}$ simply by running it simultaneously in every vertex of $\\mathbb{Z}^d$.\nWith probability $1$ this produces a labeling of $\\mathbb{Z}^d$ and it is easy to see that both conditions from the definition of ffiid are satisfied.\nOn the other hand, if $F$ is ffiid, then we collect the finite neighborhoods that determine the value at ${\\bf 0}$ together with the decision about the label.\nThis defines a pair $\\mathcal{A}_F=(\\mathcal{C},f)$ that satisfies the definition of an oblivious local algorithm.\nIt follows that the coding radius $R_\\mathcal{A}$, resp $R_F$, is the same after the translations.\nWe chose to work with the language of distributed algorithms as we borrow the techniques mainly from that area. \n\n\\begin{remark}\nIf in (b), in the definition of ffiid, we only demand that the value at ${\\bf 0}$ depends measurably on the labeling, i.e., $F$ is allowed to see the labels at the whole grid, then we say that $F$ is a \\emph{factor of iid labels (fiid)}.\n\nGiven an LCL $\\Pi$, we say that $\\Pi$ is in the class $\\mathsf{fiid}$, if there is a fiid $F$ that produces element of $X_\\Pi$ almost surely.\nIt is easy to see that $\\mathsf{TAIL}\\subseteq \\mathsf{fiid}$.\nWe remark that in the setting without inputs we have $\\mathsf{TAIL}=\\mathsf{fiid}$ on $\\mathbb{Z}$ and it is unknown, however, still plausible, whether this holds for $d>1$.\nFor general LCLs with inputs (where the definition of $\\mathsf{fiid}$ is extended appropriately) we have $\\mathsf{TAIL}\\not=\\mathsf{fiid}$ even for $d=1$, see \\cite{grebik_rozhon2021LCL_on_paths}.\n\\end{remark}\n\n\\begin{comment}\nGiven an LCL $\\Pi=(b,t,\\mathcal{P})$, we define $X_\\Pi\\subseteq b^{\\mathbb{Z}^d}$ to be the space of all $\\Pi$-colorings.\nWe say that an LCL $\\Pi$ is \\emph{satisfiable} if $\\emptyset\\not=X_\\Pi$, i.e., if there is a $\\Pi$-coloring of $\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)$.\nWe say that an LCL $\\Pi$ is \\emph{trivial} if there is $\\alpha\\in b$ such that the constant function $\\alpha$ is a $\\Pi$-coloring.\n\nThere is a canonical (right-)shift action $\\cdot$ of $\\mathbb{Z}^d$ on $b^{\\mathbb{Z}^d}$ (and in general on any $X^{\\mathbb{Z}^d}$, where $X$ is a set) that is defined as\n$$(g\\cdot x)(h)=x(h + g)$$\nfor every $g,h\\in \\mathbb{Z}^d$ and $x\\in b^{\\mathbb{Z}^d}$.\nIt is easy to see that $X_\\Pi$ is invariant under this action, i.e., $x\\in X_\\Pi$ implies $g\\cdot x\\in X_\\Pi$ for every $g\\in \\mathbb{Z}^d$.\nMoreover, since $b$ is a finite set, the product topology turns $b^\\Gamma$ into a compact metric space and it is easy to see $X_\\Pi$ is closed subspace.\nThe set $X_\\Pi$ together with the shift action is called a \\emph{sub-shift of finite type}.\n\nSuppose that $\\mathbb{Z}^d$ has an action $\\star_Y$ and $\\star_Z$ on spaces $Y$ and $Z$, respectively. \nWe say that a function $f:Y\\to Z$ is \\emph{equivariant} if $f(g\\star_Y x)=g\\star_Z f(x)$ for every $g\\in \\mathbb{Z}^d$ and $x\\in Y$.\n\nGiven an LCL $\\Pi$, our aim is to study complexity of finitary factors of iid that produce elements in $X_\\Pi$ almost surely.\nWe say that $\\Pi$ admits a \\emph{fiid solution} if there is a measurable equivariant map $F:[0,1]^{\\mathbb{Z}^d}\\to X_\\Pi$ that is defined on a $\\mu$-conull set, where $\\mu$ is the product measure on $[0,1]^{\\mathbb{Z}^d}$.\n\nWrite $\\varphi_F:[0,1]^{\\mathbb{Z}^d}\\to b$ for the composition of $F$ with the projection to ${\\bf 0}_d$-coordinate, i.e., \n$$x\\mapsto \\varphi_F(x)=F(x)({\\bf 0}_d).$$\nWe say that $F$ is a \\emph{finitary fiid (ffiid)} if for $\\mu$-almost every $x\\in [0,1]^{\\mathbb{Z}^d}$ the value of $\\varphi_F$ depends only on values of $x$ on some finite neighborhood of ${\\bf 0}_d$.\nThat is to say, there is a random variable $R_F$ that takes values in $\\mathbb{N}$, is defined $\\mu$-almost everywhere and satisfies for $\\mu$-almost every $x\\in [0,1]^{\\mathbb{Z}^d}$ the following: if $y\\in [0,1]^{\\mathbb{Z}^d}$ is such that\n$$x\\upharpoonright \\mathcal{B}_{\\operatorname{Cay}\\left(\\mathbb{Z}^d\\right)}\\left({\\bf 0}_d,R_\\varphi(x)\\right)=y\\upharpoonright \\mathcal{B}\\left({\\bf 0}_d,R_\\varphi(x)\\right),$$\nthen $\\varphi_F(x)=\\varphi_F(y)$.\n\\todo{nejak to zakoncit tim ze to nebudeme potrebovat}\n\\end{comment}\n\n\n\n\n\\begin{comment}\n\nWarm-up: Equivalence of $\\mathsf{LOCAL}(\\log^* n)$, $\\mathsf{TAIL}(\\mathrm{tower})$, $\\mathsf{CONTINUOUS}$\n\nRecall that $\\mathsf{LOCAL}(O(\\log^* n))$ denotes the class of satisfiable LCLs that satisfy (1) or (2) in Theorem~\\ref{thm:classification_of_local_problems_on_grids}.\nIn this section we show that every problem from $\\mathsf{LOCAL}(O(\\log^* n))$ is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$.\nIn the proof we use the fact that $\\mathsf{LOCAL}=\\mathsf{C-LOCAL}(O(1))$, see~Theorem~\\ref{thm:classification_of_local_problems_on_grids}.\n\nWe let $\\Delta_{s}:=|\\mathcal{B}({\\bf 0}_d,s)|$, where $s\\in \\mathbb{N}$.\nDenote as $\\mathcal{M}_s$ the subspace of $(\\Delta_s+1)^{\\mathbb{Z}^d}$ that consists of all proper vertex $(\\Delta_s+1)$-colorings of the $s$-power graph of $\\mathbb{Z}^d$.\nIt is easy to see that $\\mathcal{M}_s$ is nonempty, $\\mathbb{Z}^d$-invariant and compact.\nIn fact it is a subshift of finite type.\nThe following result holds for any graph of bounded degree, as was proved in \\cite{HolroydSchrammWilson2017FinitaryColoring}.\n\n\\begin{theorem}\\cite[Theorem~6]{HolroydSchrammWilson2017FinitaryColoring}\\label{th:ColorTower}\nLet $s\\in\\mathbb{N}$.\nThen there is an ffiid $F:[0,1]^{\\mathbb{Z}^d}\\to \\mathcal{M}_s$ and a constant $C_s>0$ such that\n$$\\mu(R_F>r)<\\frac{1}{tower(C_sr)}.$$\nIn particular, the LCL of coloring the $s$-power graph of $\\mathbb{Z}^d$ with $(\\Delta_s+1)$-many colors is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$.\n\\end{theorem}\n\n\\begin{proposition}\\label{pr:LOCALandCOLORING}\nLet $\\Pi=(b,t,\\mathcal{P})$ be a satisfiable LCL problem that is in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\nThen there is a continuous equivariant function $\\varphi:\\mathcal{M}_s\\to \\mathcal{C}_{\\Pi}$ for some $s\\in \\mathbb{N}$.\n\\end{proposition}\n\\begin{proof}\nSuppose that $\\Pi$ is in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\nBy Theorem~\\ref{thm:classification_of_local_problems_on_grids}, we have that $\\Pi$ is in the class $\\mathsf{C-LOCAL}(O(1))$.\nFix $r,\\ell\\in \\mathbb{N}$ and $\\mathcal{A}$ that witnesses this fact.\nLet $x\\in \\mathcal{M}_{r+t}$.\nThen define\n$$\\varphi(x)(h):=\\mathcal{A}([\\mathcal{B}_{\\mathbb{Z}^d_x}(h,\\ell),h]).$$\nIt is clear that $\\varphi$ is well defined, because $\\ellr)<\\frac{\\Delta_{\\ell}}{tower(C_s(r-\\ell))}\\le \\frac{1}{tower(Dr)}$$\nfor $r\\in \\mathbb{N}$ large enough and some constant $D\\in (0,+\\infty)$.\n\\end{proof}\n\n\n\n\n\ngives that\n\n\n\n\n\n\n\n\n\n$\\mathsf{TAIL}(\\tow(\\Omega(r)))$ $\\Rightarrow$ $\\mathsf{LOCAL}(O(\\log^* n))$\n\n\nFirst we show how to get randomized local algorithm from ffiid in a qualitative way.\nWe note that similar results holds for any bounded degree graph.\n\n\\begin{theorem}\\label{th:TailToRandom}\nLet $\\Pi$ be an LCL problem that is in the class $\\mathsf{TAIL}(f)$ for some $f:\\mathbb{N}\\to \\mathbb{N}$.\nSet $g(n)=\\min\\{m\\in \\mathbb{N}:f(m)\\ge n^2\\}$ and suppose that $g\\in o(n^{\\frac{1}{d}})$.\nThen $\\Pi$ is in the class $\\mathsf{R-LOCAL}(O(g(n)))$.\n\\end{theorem}\n\\begin{proof}\nFix an ffiid $F:[0,1]^{\\mathbb{Z}^d}\\to X_\\Pi$ such that\n$$\\mu(R_F>r)\\le \\frac{1}{f(r)}$$\nfor large enough $r\\in \\mathbb{N}$.\nDenote as $\\varphi_F$ the finitary rule and as $\\operatorname{dom}(\\varphi_F)$ its domain.\nThat is $\\alpha\\in \\operatorname{dom}(\\varphi_F)$ if $\\alpha\\in [0,1]^{V(\\mathcal{B}({\\bf 0}_d,s))}$ for some $s\\in \\mathbb{N}$ and satisfies\n$$F(x)=F(y)$$\nfor every extensions $\\alpha\\sqsubseteq x,y\\in [0,1]^{\\mathbb{Z}^d}$.\nIn that case we set $\\varphi_F(\\alpha)=F(x)$ for any such extension $\\alpha\\sqsubseteq x$.\n\nDefine $\\ell_n:=g(n)$.\nLet $(G,v)\\in \\operatorname{\\bf FG}_{\\bullet}^{S}$ be of a radius at most $\\ell_n$.\nConsider any map $\\ell:V(G)\\to [0,1]$.\nDefine\n$$\\mathbb{A}_n([G_l,v])=\\varphi_F(\\ell)$$\nin the case when $\\ell\\in \\operatorname{dom}(\\varphi_F)$.\nOtherwise let $\\mathbb{A}_n([G_l,v])$ to be any element in $b$.\nIt is easy to verify that this defines a randomized local algorithm.\n\nLet $G\\in \\operatorname{\\bf FG}^{S}$ be such that $|V(G)|=n$ and\n$$[\\mathcal{B}_{G}(v,\\ell_n+t),v]=\\left[\\mathcal{B}({\\bf 0}_d,\\ell_n+t),{\\bf 0}_d\\right]$$\nholds for every $v\\in V(G)$.\nNote that by the assumption on $g$ such graphs do exist, e.g., large enough torus.\nFor every $v\\in V(G)$ define the event\n$$A_v=\\left\\{\\ell\\in [0,1]^{V(G)}:\\ell\\upharpoonright \\mathcal{B}_G(v,\\ell_n)\\not \\in \\operatorname{dom}(\\varphi_F)\\right\\}.$$\nNote that if $\\ell\\not \\in \\bigcup_{v\\in V(G)} A_v$, then $\\mathbb{A}(G_l)$ is a $\\Pi$-coloring.\nTo see this pick $v\\in V(G)$ and any extension $x\\in [0,1]^{\\mathbb{Z}^d}$ of $\\ell\\upharpoonright \\mathcal{B}_G(v,\\ell_n+t)$ that is in the domain of $F$, where we interpret $v$ as the origin ${\\bf 0}_d$.\nThen $F(x)\\in X_\\Pi$ and\n$$F(x)\\upharpoonright \\mathcal{B}({\\bf 0}_d,t)\\simeq\\mathbb{A}(G_l)\\upharpoonright \\mathcal{B}_G(v,t).$$\nConsequently, we have $\\mathcal{P}([\\mathcal{B}_{G_{\\mathbb{A}(G_l)}}(v,t),v])=1$.\n\nLet $v\\in V(G)$.\nWe have $\\ell\\in A_v$ if and only if $R_F(x)>\\ell_n$ for every $x\\in [0,1]^\\Gamma$ that extends $\\ell$, i.e., every $x\\in [0,1]^{\\mathbb{Z}^d}$ that satisfies\n$$[B_{G_l}(v,\\ell_n),v]=[B_{\\mathbb{Z}^d_x}({\\bf 0}_d,\\ell_n),{\\bf 0}_d].$$ \nThis implies that\n$$\\lambda^{\\oplus n}(A_v)=\\mu\\left(\\left\\{x\\in [0,1]^{\\mathbb{Z}^d}:R_F(x)>\\ell_n\\right\\}\\right)\\le \\frac{1}{f(\\ell_n)}\\le \\frac{1}{n^2}.$$\nA union bound then shows that $\\lambda^{\\oplus n}\\left(\\bigcup_{v\\in V(G)}A_v\\right)\\le \\frac{1}{n}$ and the proof is finished.\n\\end{proof}\n\n\n\n\\begin{corollary}\\label{cor:tower=local}\n$\\mathsf{TAIL}(\\tow(\\Omega(r)))=\\mathsf{LOCAL}(O(\\log^* n))$.\n\\end{corollary}\n\\begin{proof}\nIt follows from Corollary~\\ref{cor:locallocal} itself has the following implications: \n\n\\begin{corollary}\n\\label{cor:block_factor}\n$\\Pi$ can be solved by a block factor (recall \\cref{def:tail}) if and only if $\\Pi \\in \\mathsf{LOCAL}(O(1))$. \n\\end{corollary}\n\\begin{proof}\nThis follows from \\cref{thm:classification_of_local_problems_on_grids,rem:chang_pettie_decomposition}. \n\\end{proof}\n\n\\begin{corollary}\n\\label{cor:BelowTOWERtrivial}\nIf $\\Pi \\in \\mathsf{TAIL}(\\mathrm{tower}(\\omega(r)))$, then $\\Pi \\in \\mathsf{LOCAL}(O(1))$ and if $\\Sigma_{in} = \\emptyset$ then $\\Pi$ is trivial. \n\\end{corollary}\n\\begin{proof}\n\\cref{lem:tail->local} gives that $\\Pi \\in\\mathsf{LOCAL}(o(\\log^* n))$ and \\cref{thm:classification_of_local_problems_on_grids} then implies the rest. \n\\end{proof}\n\nIn particular, Corollary~\\ref{cor:BelowTOWERtrivial} answers in negative Problem~(ii) in \\cite{HolroydSchrammWilson2017FinitaryColoring} that asked about the existence of a nontrivial local problem (without input labels) with super-tower function tail. \n\nThe second corollary is the following:\n\n\\begin{corollary}\n\\label{cor:local=tower}\n$\\Pi \\in \\mathsf{TAIL}(\\tow(\\Omega(r)))$ if and only if $\\Pi \\in \\mathsf{LOCAL}(O(\\log^* n))$. \n\\end{corollary}\n\\begin{proof}\nIf $\\Pi \\in \\mathsf{TAIL}(\\tow(\\Omega(r)))$, \\cref{lem:tail->local} yields $\\Pi \\in \\mathsf{RLOCAL}(O(\\log^* n))$. The other direction (special case of \\cref{thm:tail=uniform}) follows from \\cref{rem:chang_pettie_decomposition} that essentially states that it suffices to solve the independent set problem in the power graph. But this follows from a version of Linial's algorithm from \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n\\end{proof}\n\nFinally, we get the following. \n\\begin{corollary}\n\\label{cor:Below_d_local}\nIf $\\Pi \\in \\mathsf{TAIL}(\\omega(r^{2d}))$, then $\\Pi \\in \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$. \n\\end{corollary}\n\\begin{proof}\n\\cref{lem:tail->local} gives that $\\Pi \\in \\mathsf{LOCAL}(o(n^{1\/d}))$ and \\cref{thm:classification_of_local_problems_on_grids} then implies that $\\Pi \\in \\mathsf{LOCAL}(O(\\log^* n))$ which is equivalent to $\\Pi \\in \\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r)))$ by \\cref{cor:local=tower}. \n\\end{proof}\n\nCorollary~\\ref{cor:Below_d_local} answers in negative Problem~(iii) in \\cite{HolroydSchrammWilson2017FinitaryColoring} that asked about the existence of a problem with all its moments finite but with sub-tower function decay. Recall that \\cref{fact:moment_implies_tail} implies that if the $2d$-th moment is finite, then the tail is $\\omega(r^{2d})$ and, hence, the problem can in fact be solved with a tower tail decay. \nHowever, the $2d$ is not the right dependency and in the next section we prove a tight version of this result.\n\\subsection{Tight Grid Speedup for $\\mathsf{TAIL}$}\n\\label{subsec:tail_d-1_speedup}\n\nWhile \\cref{cor:Below_d_local} already answers the Problem of \\cite{HolroydSchrammWilson2017FinitaryColoring}, we now improve the speedup from the tail of $\\omega(r^{d-1})$. This result is tight (up to $o(1)$ in the exponent, see the discussion in \\cref{subsec:toast_construction} and \\cref{cor:toast_lower_bound}) and it is the first nontrivial difference between the $\\mathsf{TAIL}$ and $\\mathsf{LOCAL}$ model that we encountered so far. \n\n\n\n\n\\begin{theorem}\n\\label{thm:grid_speedup}\nLet $d,t\\in \\mathbb{N}$.\nFor $d \\ge 2$ there is $C_{d,t}>0$ such that if an LCL $\\Pi$ is in $\\mathsf{TAIL}(C_{d,t}r^{d-1})$, then $\\Pi \\in \\mathsf{TAIL}(\\tow(\\Omega(r)))$.\n\nSimilarly, for $d=1$ we have that a local coloring problem $\\Pi$ admits a solution by an oblivious random algorithm (it is in $\\mathsf{TAIL}(f(r))$ for any function $f = \\omega(1)$) if and only if $\\Pi$ is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$.\n\\end{theorem}\nNote that \\cref{fact:moment_implies_tail} hence shows that ${\\bf M}(\\Pi) > d-1$ already implies that $\\Pi$ is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$. \n\nThe proof is a simple adaptation of Lov\u00e1sz local lemma speedup results such as \\cite{chang2017time_hierarchy,fischer2017sublogarithmic}. \nThe proof goes as follows. Split the grid into boxes of side length roughly $r_0$ for some constant $r_0$. This can be done with $O(\\log^* (n))$ local complexity (or, equivalently, with tower tail decay). \n\nNext, consider a box $B$ of side length $r_0$ and volume $r_0^d$. In classical speedup argument we would like to argue that we can deterministically set the randomness in nodes of $B$ and its immediate neighborhood so that we can simulate the original oblivious local algorithm $\\mathcal{A}$ with that randomness and it finishes in at most $r_0$ steps. \nIn fact, ``all nodes of $B$ finish after at most $r_0$ steps'' is an event that happens with constant probability if the tail decay of $\\mathcal{A}$ is $r^d$, since then each node fails with probability $1\/r_0^d$ and there are $r_0^d$ nodes in $B$. \n\nHowever, we observe that it is enough to consider only the boundary of $B$ with $O(r_0^{d-1})$ nodes and argue as above only for the boundary nodes, hence the speedup from $r^{d-1}$ tail decay. \nThe reason for this is that although the coding radius of $\\mathcal{A}$ inside $B$ may be very large, just the fact that there \\emph{exists} a way of filling in the solution to the inside of $B$, given the solution on its boundary, is enough for our purposes. The new algorithm simply simulates $\\mathcal{A}$ only on the boundaries and then fills in the rest. \n\n\\begin{proof}\n\n\nWe will prove the result for $d \\ge 2$, the case $d=1$ is completely analogous and, in fact, simpler. \nSet $C=C_{d,t}:=d\\cdot t\\cdot(2\\Delta)^{2\\Delta}$, where $\\Delta=20^d$.\nSuppose that there is an oblivious local algorithm $\\mathcal{A}$ that explores the neighborhood of a given vertex $u$ and outputs a color of $u$ after seeing a radius $R$ such that \\begin{equation}\\label{eq:tail}\n\\P(R > r)\\le \\frac{1}{C r^{d-1}}\n\\end{equation}\nfor $r\\in \\mathbb{N}$ large enough.\nLet $r_0\\in \\mathbb{N}$ be large enough such that \\eqref{eq:tail} is satisfied.\nWrite $\\Box_{r_0}$ for the local coloring problem of tiling the input graph (which is either a large enough torus or an infinite grid $\\mathbb{Z}^d$) into boxes with length sizes in the set $\\{r_0,r_0+1\\}$ such that, additionally, each box $B$ is colored with a color from $[\\Delta+1]$ which is unique among all boxes $B'$ with distance at most $3(r_0+t)$. \n\nThe problem of tiling with boxes of side lengths in $\\{r_0, r_0+1\\}$ can be solved in $O(\\log^* n)$ local complexity \\cite[Theorem~3.1]{GaoJackson} but we only sketch the reason here: Find a maximal independent set of a large enough power of the input graph $\\mathbb{Z}^d$. Then, each node $u$ in the MIS collects a Voronoi cell of all non-MIS nodes such that $u$ is their closest MIS node. As the shape of the Voronoi cell of $u$ depends only on other MIS vertices in a constant distance from $u$, and their number can be bounded as a function of $d$ (see the ``volume'' argument given next), one can discretize each Voronoi cell into a number of large enough rectangular boxes. Finally, every rectangular box with all side lengths of length at least $r^2$ can be decomposed into rectangular boxes with all side lengths either $r$ or $r+1$. \n\nNote that in additional $O(\\log^* n)$ local steps, the boxes can then be colored with $(\\Delta+1)$-many colors in such a way that boxes with distance at most $3(r_0+t)$ must have different colors.\nTo see this, we need to argue that there are at most $\\Delta$ many boxes at distance at most $3(r_0+t)$ from a given box $B$.\nThis holds because if $B,B'$ are two boxes and $d(B,B')\\le 3(r_0+t)$, then $B'$ is contained in the cube centered around the center of $B$ of side length $10(r_0+t)$.\nThen a volume argument shows that there are at most $\\Delta\\ge \\frac{(10(r_0+t))^d}{r_0^d}$ many such boxes $B'$ (when $B$ is fixed).\n\nWe show now that there is a local algorithm $\\mathcal{A'}$ that solves $\\Pi$ when given the input graph $G$ together with a solution of $\\Box_{r_0}$ as an input.\nThis clearly implies that $\\Pi$ is in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\n\n\n\nLet $\\{B_i\\}_{i\\in I}$ be the input solution of $\\Box_{r_0}$, together with the coloring $B_i\\mapsto c(B_i)\\in [\\Delta+1]$.\nWrite $V(B)$ for the vertices of $B$ and $\\partial B$ for the vertices of $B$ that have distance at most $t$ from the boundary of $B$.\nWe have $|\\partial B|< d\\cdot t \\cdot (r_0+1)^{d-1}$.\nDefine $N(B)$ to be the set of rectangles of distance at most $(r_0+t)$ from $B$.\n\n\\begin{figure}\n \\centering\n \\includegraphics[width= .9\\textwidth]{fig\/tail_LLL.eps}\n \\caption{The picture shows one block $B$ with a red boundary $\\partial B$ of length $t$. Boxes in $N(B)$, that is with distance at most $(r_0+t)$ from $B$ are dark grey. }\n \\label{fig:tail_LLL}\n\\end{figure}\n\nWe are now going to describe how $\\mathcal{A}'$ iterates over the boxes in the order of their colors and for each node $u$ in the box $B$ it carefully fixes a random string $\\mathfrak{r}(u)$. \nThis is done in such a way that after fixing the randomness in all nodes, we can apply the original algorithm $\\mathcal{A}$ to each vertex in $\\partial B$ for each $B$ with the randomness given by $\\mathfrak{r}$ and the local algorithm at $u \\in \\partial B$ will not need to see more than its $r$-hop neighborhood. \n\nGiven a partial function $\\mathfrak{r_\\ell}:V(G)\\to [0,1]$ corresponding to nodes such that their randomness was already fixed, every box $B$ can compute a probability $\\mathbb{P}(\\fE_{\\ell}(B))$, where $\\fE_{\\ell}(B)$ is the (bad) event that there is a vertex in $\\partial B$ that needs to look further than $r_0$ to apply the local algorithm $\\mathcal{A}$.\nNote that for computing $\\mathbb{P}(\\fE_{\\ell}(B))$ it is enough to know the $r_0$-neighborhood of $B$.\n\nThe algorithm $\\mathcal{A}$ has $\\Delta+1$ rounds.\nIn the $0$-th round, we set $\\mathfrak{r}_0$ to be the empty function.\nNote that the probability that a vertex needs to look further than $r_0$ in the algorithm $\\mathcal{A}$ is smaller than $\\frac{1}{C r_0^{d-1}}$ by \\eqref{eq:tail}.\nThe union bound gives $\\mathbb{P}(\\fE_0(B))<\\frac{1}{(2\\Delta)^{2\\Delta}}$ for every $B$.\n\nIn each step $\\ell\\in [\\Delta+1]$ we define $\\mathfrak{r}_B:V(B)\\to [0,1]$ for every $B$ such that $c(B)=\\ell$ and put $\\mathfrak{r}_l$ to be the union of all $\\mathfrak{r}_D$ that have been defined so far, i.e., for $D$ such that $c(D)\\le \\ell$.\nWe make sure that the following is satisfied for every $B$:\n\\begin{enumerate}\n\t\\item [(a)] the assignment $\\mathfrak{r}_B$ depends only on $2(r_0+t)$-neighborhood of $B$ (together with previously defined $\\mathfrak{r}_{\\ell-1}$ on that neighborhood) whenever $c(B)=\\ell$,\n\t\\item [(b)] $\\mathbb{P}(\\fE_l(B))<\\frac{1}{(2\\Delta)^{2\\Delta-\\ell}}$,\n\t\\item [(c)] restriction of $\\mathfrak{r}_l$ to $N(B)$ can be extend to the whole $\\mathbb{Z}^d$ such that $\\mathcal{A}$ is defined at every vertex, where we view $N(B)$ embedded into $\\mathbb{Z}^d$.\n\\end{enumerate}\nWe show that this gives the required algorithm.\nIt is easy to see that by (a) this defines a local construction that runs in constantly many rounds and produces a function $\\mathfrak{r}:=\\mathfrak{r}_{\\Delta+1}:V(G)\\to [0,1]$.\nBy (b) every vertex in $\\partial B$ for some $B$ looks into its $r_0$ neighborhood and use $\\mathcal{A}$ to choose its color according to $\\mathfrak{r}$.\nFinally, by (c) it is possible to assign colors in $B\\setminus \\partial B$ for every $B$ to satisfy the condition of $\\Pi$ within each $B$.\nPicking minimal such solution for each $B\\setminus \\partial B$ finishes the description of the algorithm.\n\nWe show how to proceed from $0\\le \\ell<\\Delta+1$ to $\\ell+1$.\nSuppose that $\\mathfrak{r}_l$ satisfies the conditions above.\nPick $B$ such that $c(B)=\\ell+1$.\nNote that $\\mathfrak{r}_B$ that we want to find can influence rectangles $D\\in N(B)$ and these are influenced by the function $\\mathfrak{r}_l$ at distance at most $r_0$ from them.\nAlso given any $D$ there is at most one $B\\in N(D)$ such that $c(B)=\\ell+1$. \nWe make the decision $\\mathfrak{r}_B$ based on $2r_0$ neighborhood around $B$ and the restriction of $\\mathfrak{r}_l$ there, this guarantees that condition (a) is satisfied.\n\nFix $D\\in N(B)$ and pick $\\mathfrak{r}_B$ uniformly at random.\nWrite $\\mathfrak{S}_{B,D}$ for the random variable that computes the probability of $\\fE_{\\mathfrak{r}_B\\cup \\mathfrak{r}_l}(D)$.\nWe have $\\mathbb{E}(\\mathfrak{S}_{B,D})=\\mathbb{P}(\\fE_l(D))<\\frac{1}{(2\\Delta)^{2\\Delta-\\ell}}$.\nBy Markov inequality we have\n$$\\mathbb{P}\\left(\\mathfrak{S}_{B,D}\\ge 2\\Delta \\frac{1}{(2\\Delta)^{2\\Delta-\\ell}}\\right)<\\frac{1}{2\\Delta}.$$\nSince $|N(B)|\\le \\Delta$ we have that with non-zero probability $\\mathfrak{r}_B$ satisfies $\\mathfrak{S}_{B,D}<\\frac{1}{(2\\Delta)^{2\\Delta-\\ell-1}}$ for every $D\\in N(B)$.\nSince $\\mathcal{A}$ is defined almost everywhere we find $\\mathfrak{r}_B$ that satisfies (b) and (c) above.\nThis finishes the proof. \n\\end{proof}\n\nAs we already mentioned, the proof is in fact an instance of the famous Lov\u00e1sz Local Lemma problem. This problem is extensively studied in the distributed computing literature \\cite{chang2016exp_separation, fischer2017sublogarithmic, brandt_maus_uitto2019tightLLL, brandt_grunau_rozhon2020tightLLL} as well as descriptive combinatorics \\cite{KunLLL,OLegLLL,MarksLLL,Bernshteyn2021LLL}.\nOur usage corresponds to a simple LLL algorithm for the so-called exponential criteria setup (see \\cite{fischer2017sublogarithmic} and \\cite{Bernshteyn2021local=cont}). \n\n\n \n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{$\\TOAST$: General Results}\n\\label{sec:Toast}\n\nIn this section we introduce the toast construction. This construction is very standard in descriptive combinatorics \\cite{FirstToast,GJKS,Circle} and the theory of random processes (there it is often called \\emph{the hierarchy}) \\cite{HolroydSchrammWilson2017FinitaryColoring,ArankaFeriLaci,Spinka}.\n\nWhile the construction itself is well-known, in this paper we study it as a class of problems by itself. An analogy to our work is perhaps the paper of \\cite{ghaffari2018derandomizing} that defines the class $\\mathsf{S-LOCAL}$ and shows how to derandomize it.\n\n\n\nInformally, the basic toast construction works as follows. There is an adversary that has access to the (infinite) input grid $\\mathbb{Z}^d$. The adversary is sequentially choosing some connecting pieces of the grid and gives them to our algorithm $\\mathcal{A}$. The algorithm $\\mathcal{A}$ needs to label all vertices in the piece $D$ once it is given $D$. \nNote that $\\mathcal{A}$ does not know where exactly $D$ lies in $\\mathbb{Z}^d$ and its decision cannot depend on the precise position of $D$ in the input graph. In general, the piece $D$ can be a superset of some smaller pieces that were already colored by $\\mathcal{A}$. In this case, $\\mathcal{A}$ needs to extend their coloring to the whole $D$. This game is played for countable number of steps until the whole grid is colored. The algorithm solves $\\Pi$ if it never happens during the procedure that it cannot extend the labeling to $D$ in a valid manner. \n\n\n\n\n\nThe construction sketched above is what we call a spontaneous deterministic toast or just $\\TOAST$ algorithm. \nIn the next subsection we give a formal definition of $\\TOAST$ algorithms and we follow by proving relations between $\\TOAST$ algorithms and $\\mathsf{TAIL}$ complexities. \n\n\n\\subsection{Formal Definition of $\\TOAST$ classes}\n\\label{subsec:toast_definition}\nWe now give a formal definition of our four $\\TOAST$ classes. First we define the hereditary structure on which the $\\TOAST$ algorithm is interpreted. This is also called a toast (we use lowercase letters) and is parameterized by the minimum distance of boundaries of different pieces of the toast. \nThis structure also corresponds to the adversary that asks our algorithm to extend a partial solution to $\\Pi$ to bigger and bigger pieces used in the informal introduction. \n\n\\begin{definition}[$q$-toast on $\\mathbb{Z}^d$]\nFor a $q >0$, a \\emph{$q$-toast on $\\mathbb{Z}^d$} is a collection $\\mathcal{D}$ of finite connected subsets of $\\mathbb{Z}^d$ that we call \\emph{pieces} such that\n\\begin{itemize}\n \\item for every $g,h\\in \\mathbb{Z}^d$ there is a piece $D\\in\\mathcal{D}$ such that $g,h\\in D$,\n \\item if $D,E\\in \\mathcal{D}$ and $d(\\partial D,\\partial E)\\le q$, then $D=E$, where $\\partial D$ is a boundary of a set $D$.\n\\end{itemize}\n\\end{definition}\n\nGiven a $q$-toast $\\mathcal{D}$, the \\emph{rank} of its piece $D \\in \\mathcal{D}$ is the level in the toast hierarchy that $D$ belongs to in $\\mathcal{D}$. Formally, we define a \\emph{rank function $\\mathfrak{r}_\\mathcal{D}$ for the toast $\\mathcal{D}$} as $\\mathfrak{r}_\\mathcal{D}:\\mathcal{D}\\to \\mathbb{N}$, inductively as\n$$\\mathfrak{r}_\\mathcal{D}(D)=\\max\\{\\mathfrak{r}_{\\mathcal{D}}(E)+1:E\\subseteq D, \\ E\\in \\mathcal{D}\\}$$ \nfor every $D\\in \\mathcal{D}$.\n\nNext, $\\TOAST$ algorithms are described in the following language of extending functions. \n\n\\begin{definition}[$\\TOAST$ algorithm]\n\\label{def:extending_function}\nWe say that a function $\\mathcal{A}$ is a \\emph{$\\TOAST$ algorithm} for a local problem $\\Pi = (\\Sigma_{in},\\Sigma_{out}, t, \\mathcal{P})$ if \n\\begin{enumerate}\n \\item [(a)] the domain of $\\mathcal{A}$ consists of pairs $(A, c)$, where $A$ is a finite connected $\\Sigma_{in}$-labeled subset of $\\mathbb{Z}^d$ and $c$ is a partial map from $A$ to $\\Sigma_{out}$,\n \\item [(b)] $\\mathcal{A}(A,c)$ is a partial coloring from $A$ to $\\Sigma_{out}$ that extends $c$, i.e., $\\mathcal{A}(A,c)\\upharpoonright \\operatorname{dom}(c)=c$, and that depends only on the isomorphism type of $(A,c)$ (not on the position in $\\mathbb{Z}^d$). \n \n\\end{enumerate}\nA \\emph{randomized $\\TOAST$ algorithm} is moreover given as an additional input a labeling of $A$ with real numbers from $[0,1]$ (corresponding to a stream of random bits). \nFinally, if $\\mathcal{A}$ has the property that $\\operatorname{dom}(\\mathcal{A}(A,c)) =A$ for every pair $(A,c)$, then we say that $\\mathcal{A}$ is a \\emph{spontaneous (randomized) $\\TOAST$ algorithm}.\n\\end{definition}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=\\textwidth]{fig\/3-col.eps}\n \\caption{A $q$-toast and an example of an extending function for $3$-coloring. Observe that our extending function does not give a correct algorithm for $3$-coloring since the piece on the right is contained in an even bigger piece and the $3$-coloring constructed on the right cannot be extended to the node marked with an arrow}\n \\label{fig:toast_example}\n\\end{figure}\n\n\nIt is easy to see that a $\\TOAST$ algorithm $\\mathcal{A}$ can be used to color ($\\Sigma_{in}$-labeled) $\\mathbb{Z}^d$ inductively on the rank $\\mathfrak{r}_\\mathcal{D}$ of any $q$-toast $\\mathcal{D}$ for every $q\\in \\mathbb{N}$. \nNamely, define $\\mathcal{A}(\\mathcal{D})$ to be the partial coloring of $\\mathbb{Z}^d$ that is the union of $\\mathcal{A}(D)$, where $D\\in \\mathcal{D}$ and $\\mathcal{A}(D)$ is defined inductively on $\\mathfrak{r}_\\mathcal{D}$.\nFirst put $\\mathcal{A}(D):=\\mathcal{A}(D,\\emptyset)$ whenever $\\mathfrak{r}_\\mathcal{D}(D)=0$.\nHaving defined $\\mathcal{A}(D)$ for every $D\\in \\mathcal{D}$ such that $\\mathfrak{r}_\\mathcal{D}(D)=k$, we define\n$$\\mathcal{A}(D'):=\\mathcal{A}\\left(\\bigcup\\{E\\in \\mathcal{D}:E\\subseteq D'\\},\\bigcup\\{\\mathcal{A}(E):E\\subseteq D', \\ E\\in \\mathcal{D}\\}\\right).$$\nfor every $D'\\in \\mathcal{D}$ such that $\\mathfrak{r}_\\mathcal{D}(D')=k+1$. That is, we let $\\mathcal{A}$ to extend the partial solution to $\\Pi$ to bigger and bigger pieces of the toast $\\mathcal{D}$. \nObserve that if $\\mathcal{A}$ is a spontaneous $\\TOAST$ function, then $\\operatorname{dom}(\\mathcal{A}(\\mathcal{D}))=\\mathbb{Z}^d$ for every $q$-toast $\\mathcal{D}$ and every $q\\in \\mathbb{N}$.\n\n\\paragraph{$\\TOAST$ Algorithms}\nWe are now ready to define the four types of toast constructions that we are going to use (cf. \\cref{fig:big_picture_grids}). \n\n\\begin{definition}[$\\LTOAST$]\nLet $\\Pi$ be an LCL.\nWe say that $\\Pi$ is in the class $\\LTOAST$ if there is $q>0$ and an $\\TOAST$ algorithm $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is a $\\Pi$-coloring for every $q$-toast $\\mathcal{D}$.\n\nMoreover, we say that $\\Pi$ is in the class $\\TOAST$ if there is $q\\in \\mathbb{N}$ and a spontaneous $\\TOAST$ algorithm $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is a $\\Pi$-coloring for every $q$-toast $\\mathcal{D}$.\n\\end{definition}\n\nBefore we define a randomized version of these notions, we make several comments.\nFirst note that $\\TOAST\\subseteq \\LTOAST$.\nBy (c), if we construct a ($q$-)toast using an oblivious local algorithm, then we get instantly that every problem from the class $\\TOAST$ is in the class $\\mathsf{TAIL}$, see Section~\\ref{sec:Toast}.\nSimilar reasoning applies to the class $\\mathsf{BOREL}$, see Section~\\ref{sec:OtherClasses}.\nWe remark that we are not aware of any LCL that is not in the class $\\TOAST$.\n\n\\begin{comment}\n\\begin{definition}[Randomized toast algorithm]\n\\label{def:extending_function}\nLet $b$ be a set.\nWe say that a function $\\mathcal{A}$ is a randomized $\\TOAST$ algorithm if \n\\begin{enumerate}\n \\item [(a)] the domain of $\\mathcal{A}$ consists of triplets $(A, c, \\ell)$, where $A$ is a finite connected subset of $\\mathbb{Z}^d$, $c$ is a partial coloring of $A$ and $\\ell:A\\to [0,1]$. \n \\item [(b)] $\\mathcal{A}(A,c,\\ell)$ is a partial coloring of $A$ that extends $c$, i.e., $\\mathcal{A}(A,c,\\ell)\\upharpoonright \\operatorname{dom}(c)=c$. \n \\item [(c)] the function $\\mathcal{A}$ is invariant with respect to the shift action of $\\mathbb{Z}^d$ on $[0,1]^{\\mathbb{Z}^d)}$. \n\\end{enumerate}\nMoreover, if $\\mathcal{A}$ has the property that $\\operatorname{dom}(\\mathcal{A}(A,c,\\ell)) =A$ for every triplet $(A,c,\\ell)$, then we say that $\\mathcal{A}$ is a \\emph{spontaneous randomized toast algorithm}.\n\\end{definition}\n\nSimilarly as above, we define $\\mathcal{A}(\\mathcal{D},\\ell)$ for every $q$-toast $\\mathcal{D}$, $q\\in \\mathbb{N}$ and $\\ell\\in [0,1]^{\\mathbb{Z}^d}$.\nObserve that if $\\mathcal{A}$ is a spontaneous randomized toast algorithm, then $\\operatorname{dom}(\\mathcal{A}(\\mathcal{D},\\ell))=\\mathbb{Z}^d$.\n\\end{comment}\n\n\\begin{definition}[$\\RLTOAST$]\nLet $\\Pi$ be an LCL.\nWe say that $\\Pi$ is in the class $\\RLTOAST$ if there is $q>0$ and a randomized $\\TOAST$ algorithm $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is a $\\Pi$-coloring for every $q$-toast $\\mathcal{D}$ with probability $1$.\nThat is, after $\\Sigma_{in}$-labeling is fixed, with probability $1$ an assignment of real numbers from $[0,1]$ to $\\mathbb{Z}^d$ successfully produces a $\\Pi$-coloring against all possible toasts $\\mathcal{D}$ when $\\mathcal{A}$ is applied.\n\nMoreover, we say that $\\Pi$ is in the class $\\RTOAST$ if there is $q\\in \\mathbb{N}$ and a spontaneous randomized $\\TOAST$ algorithm $\\mathcal{A}$ as above.\n\\end{definition}\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{$\\mathsf{TAIL}$ Construction of $\\TOAST$}\n\\label{subsec:toast_construction_in_tail_sketch}\n\n\nThe class $\\TOAST$ (or $\\RTOAST$) is especially useful since in \\cref{sec:SpecificProblems} we discuss how many local problems that are not in class $O(\\log^* n)$ but on the other hand are not ``clearly hard'' like vertex $2$-coloring are in the class $\\TOAST$. Importantly, in \\cref{subsec:toast_construction} we show how one can construct the toast decomposition with an oblivious local algorithm in $\\mathsf{TAIL}(r^{1 - O(1\/\\sqrt{\\log r})})$, which yields the following theorem.\n\n\\begin{restatable}{theorem}{toastconstruction}\n\\label{thm:toast_in_tail}\nLet $\\Pi$ be an LCL that is in the class $\\TOAST$.\nThen \\[\n\\Pi \\in \\mathsf{TAIL}\\left(\\frac{r}{2^{O(\\sqrt{\\log r})}}\\right) = \\mathsf{TAIL}\\left(r^{1-o(1)}\\right)\n\\]\nand, hence, ${\\bf M}(\\Pi)\\ge 1$.\n\\end{restatable}\n\nWe observe in \\cref{cor:toast_lower_bound} that faster then linear tail decay is not possible. \n\n\\paragraph{Proof Sketch}\n\nHere we provide an intuition for the proof that is otherwise technical. We first check that maximal independent sets of power graphs of the grid $G = \\mathbb{Z}^d$ can be constructed effectively, e.g. because of the MIS algorithm of Ghaffari \\cite{ghaffari2016MIS}. The slowdown by this construction can be neglected in this high-level overview. \n\nWe then take a sequence $r_k = 2^{k^2}$ and construct sparser and sparser MIS in graphs $G^{r_k}, k = 1, 2, \\dots$. After constructing MIS of $G^{r_k}$, we construct Voronoi cells induced by this MIS. These Voronoi cells $V_1, V_2, \\dots$ will be, roughly, the pieces of rank $k$, but we need to ensure that all the boundaries are at least $q$ hops apart. \n\nTo do so, we first shrink each Voronoi cells by $t_k = O(\\sum_{j < k} r_j) \\ll r_k$ hops (grey boundary area in \\cref{fig:toast_construction}) and get new, smaller, cells $C_1, C_2, \\dots$. The final sequence of pieces $H_1, H_2, \\dots$ is then defined inductively, namely to define $H_i$ we start with $C_1$, then add all Voronoi cells from $k-1$-th iteration that $C_1$ intersects, then add all Voronoi cells from $k_2$-th iteration that the new, bigger, cluster intersects, \\dots\nThe definition of $t_k$ ensures that the final pieces constructed after all $k$ iterations are sufficiently far from each other, and also compatible with the smaller rank pieces. In fact, we always have $H_i \\subseteq V_i$. \n\nIt remains to argue that after $k$ steps of this process, only a small fraction of nodes remains uncovered. This is done roughly as follows. The vertices from $V_i \\setminus C_i$ that are not covered by $C_i$s are roughly $t_k \/ r_k \\approx A r_{k-1} \/ r_k$ fraction of vertices where $A$ is some (large) constant. By induction, $1 - A r_{k-2} \/ r_{k-1}$ fraction of those vertices were covered in step $k-1$ and so on. The fraction of remaining vertices is then $\\frac{A r_{k-1}}{r_{k}} \\cdot \\frac{A r_{k-2}}{r_{k-1}} \\dots \\frac{A }{r_1} = \\frac{A^k}{r_k}$. This fraction is asymptotically almost equal to $1\/r_k$, which is the desired tail behavior. However, the additional term $A^k$ is the reason why the final tail decay for $r$ such that $r_k \\le r \\le r_{k+1} = 2^{(k+1)^2}$ is bounded by $2^{O(k)}\/2^{k^2} \\le \\left( \\frac{1}{r} \\right)^{\\frac{k^2 - O(k)}{(k+1)^2}} = \\left( \\frac{1}{r} \\right)^{1 - O(1\/k)} = \\left( \\frac{1}{r} \\right)^{1 - O(1\/\\sqrt{\\log r})} $, as needed. \n\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width = .95\\textwidth]{fig\/toast_construction.eps}\n \\caption{A step during the toast construction. Voronoi cell $V$ is first shrinked to $C$ and we get the final piece $H$ by taking the union of smaller-level Voronoi cells intersecting $C$ (in this picture, there is only one lower level but the construction is recursive in general). }\n \\label{fig:toast_construction}\n\\end{figure}\n\n\n\n\n\n\n\n\n\n\n\\subsection{$\\RTOAST = \\TOAST$}\n\\label{subsec:rtoast=toast}\n\n\n\nIn this section we show that randomness does not help spontaneous toast algorithms.\n\n\\begin{restatable}{theorem}{toastrtoast}\\label{thm:derandomOfToast}\nLet $\\Pi$ be an LCL.\nThen $\\Pi$ is in the class $\\RTOAST$ if and only if $\\Pi$ is in the class $\\TOAST$.\n\\end{restatable}\n\nThe proof goes as follows. We start with a random toast algorithm $\\mathcal{A}$ and make it deterministic as follows. When the deterministic toast algorithm gets a piece $D$ of rank $1$, it considers all possible outputs that the algorithm $\\mathcal{A}$ could produce given the input random string and it uses any one of them that $\\mathcal{A}$ would produce with nonzero probability. \nBy induction, whenever $\\mathcal{A}'$ later gets a piece $D$ containing some pieces of smaller rank with a partial solution that $\\mathcal{A}'$ already produced, the original algorithm $\\mathcal{A}$ could have produced the same partial solution with nonzero probability, hence $\\mathcal{A}$ can extend the coloring in a valid manner and $\\mathcal{A}'$ deterministically choosing any such valid extension that $\\mathcal{A}$ would use with nonzero probability. \nThis shows that $\\mathcal{A}'$ also solves $\\Pi$. \nA formal proof follows. \n\n\\begin{proof}\n\nIt is easy to see that $\\TOAST\\subseteq \\RTOAST$.\nSuppose that $\\Pi$ is in the class $\\RTOAST$ and fix a witnessing $q\\in \\mathbb{N}$ together with a spontaneous randomized $\\TOAST$ algorithm $\\mathcal{A}$.\nWe may assume that $t$, the locality of $\\Pi$, satisfies $t\\le q$.\n\nLet $(A,c)$ be such that $A\\subseteq \\mathbb{Z}^d$ is finite and connected and $\\operatorname{dom}(c)\\subseteq A$.\nWe say that $(A,c)$ comes from a \\emph{partial random $q$-toast} if there is a collection $\\mathcal{D}$ of subsets of $A$ together with an event $\\mathcal{F}$ on $\\operatorname{dom}(c)$ such that\n\\begin{itemize}\n \\item if $D,E\\in \\mathcal{D}$ and $d(\\partial D,\\partial E)\\le q$, then $D=E$,\n \\item $A\\in \\mathcal{D}$,\n \\item $\\bigcup_{E\\in \\mathcal{D}\\setminus \\{A\\}}E=\\operatorname{dom}(c)$,\n \\item $\\mathcal{F}$ has non-zero probability and if we apply inductively $\\mathcal{A}$ on $\\operatorname{dom}(c)$ we obtain $c$.\n\\end{itemize}\nIf $(A,c)$ comes from a partial random $q$-toast, then we let $\\mathcal{E}(A,c)$ to be any event $\\mathcal{F}$ as above.\nOtherwise we put $\\mathcal{E}(A,c)$ to be any event on $\\operatorname{dom}(c)$ of non-zero probability.\nIt is easy to see that $\\mathcal{E}$ does not depend on the position of $A$ in $\\mathbb{Z}^d$.\nThis is because $\\mathcal{A}$ is such.\n\nWe now define a spontaneous toast algorithm $\\mathcal{A}'$.\nGiven $(A,c)$, we let $\\mathcal{F}(A,c)$ to be any event on $A$ that has non-zero probability, satisfies\n$$\\mathcal{F}(A,c)\\upharpoonright \\operatorname{dom}(c)\\subseteq \\mathcal{E}(A,c)$$\nand $\\mathcal{A}(A,c)$ gives the same outcome on $\\mathcal{F}(A,c)$.\nLet $\\mathcal{A}'(A,c)$ be equal to this outcome on $\\mathcal{F}(A,c)$.\nIt is clear that $\\mathcal{A}'$ is a spontaneous $\\TOAST$ algorithm.\nObserve that if $(A,c)$ comes from a partial random $q$-toast, then $(A,\\mathcal{A}'(A,c))$ comes from a partial random $q$-toast.\nNamely, use the same $q$-toast witness $\\mathcal{D}$ as for $(A,c)$ together with the event $\\mathcal{F}(A,c)$.\n\nIt remains to show that $\\mathcal{A}'(\\mathcal{D})$ is a $\\Pi$-coloring for every $q$-toast $\\mathcal{D}$.\nInduction on $\\mathfrak{r}_{\\mathcal{D}}$ shows that $(D,\\mathcal{A}'(D))$ comes from a partial random $q$-toast for every $D\\in \\mathcal{D}$.\nThis is easy when $\\mathfrak{r}_{\\mathcal{D}}(D)=0$ and follows from the observation above in general.\nLet $v\\in \\mathbb{Z}^d$ and pick a minimal $D\\in \\mathcal{D}$ such that there is $D'\\in \\mathcal{D}$ that satisfies $v\\in D'$ and $D'\\subseteq D$.\nSince $t\\le q$ we have that $\\mathcal{B}(v,t)\\subseteq D$.\nSince $(D,\\mathcal{A}'(D))$ comes from a partial random $q$-toast it follows that there is a non-zero event $\\mathcal{F}$ such that \n$$\\mathcal{A}'(D)\\upharpoonright \\mathcal{B}(v,t)=\\mathcal{A}(\\mathcal{D})\\upharpoonright \\mathcal{B}(v,t)$$\non $\\mathcal{F}$.\nThis shows that the constraint of $\\Pi$ at $v$ is satisfied and the proof is finished.\n\\end{proof}\n\n\n\n\n\n\n\n\\subsection{$\\mathsf{TAIL} = \\RLTOAST$}\n\\label{subsec:tail=randomlazytoast}\n\nIn this section we prove that the class $\\RLTOAST$ in fact covers the whole class $\\mathsf{TAIL} = \\bigcup_{f = \\omega(1)} \\mathsf{TAIL}(f(r))$. \n\n\\begin{restatable}{theorem}{tailrltoast}\n\\label{thm:tail=rltoast}\nLet $\\Pi$ be an LCL.\nThen $\\Pi$ is in the class $\\RLTOAST$ if and only if $\\Pi$ is in the class $\\mathsf{TAIL}$.\n\\end{restatable}\n\nWe first give an informal proof. Suppose that there is a random lazy $\\TOAST$ algorithm $\\mathcal{A}$ for $\\Pi$. We construct an oblivious algorithm with nontrivial tail behavior $\\mathcal{A}'$ for $\\Pi$ by constructing a toast via \\cref{thm:toast_in_tail} and during the construction, each node moreover simulates the algorithm $\\mathcal{A}$. The algorithm $\\mathcal{A}$ finishes after some number of steps with probability $1$, hence $\\mathcal{A}'$ finishes, too, with probability $1$ for each node. \n\nOn the other hand, suppose that there is oblivious local algorithm $\\mathcal{A}'$ for $\\Pi$. \nThen we can construct a lazy randomized toast algorithm $\\mathcal{A}$ that simply simulates $\\mathcal{A}'$ for every node. This is possible since as the pieces around a given vertex $u$ that has gets are growing bigger and bigger, the radius around $u$ that $\\mathcal{A}$ sees also grows, until it is large enough and $\\mathcal{A}$ sees the whole coding radius around $u$ and applies $\\mathcal{A}'$ then. \n\nA formal proof follows. \n\n\n\\begin{proof}\nSuppose that $\\Pi$ is in the class $\\RLTOAST$, i.e, there is $q\\in\\mathbb{N}$ together with a randomized toast algorithm $\\mathcal{A}$.\nIt follows from Theorem~\\ref{thm:toast moment} that there is an oblivious local algorithm $F$ that produces $q$-toast with probability $1$.\n\nDefine an oblivious local algorithm $\\mathcal{A}'$ as follows.\nLet $\\mathcal{B}(v,r)$ be a $\\Sigma_{in}$-labeled neighborhood of $v$.\nAssign real numbers from $[0,1]$ to $\\mathcal{B}(v,r)$ and run $F$ and $\\mathcal{A}$ on $\\mathcal{B}(v,r)$.\nLet $a_{v,r}$ be the supremum over all $\\Sigma_{in}$-labelings of $\\mathcal{B}(v,r)$ of the probability that $\\mathcal{A}'$ fails.\nIf $a_{v,r}\\to 0$ as $r \\to \\infty$, then $\\mathcal{A}'$ is an oblivious local algorithm.\n\nSuppose not, i.e., $a_{v,r}>a>0$.\nSince $\\Sigma_{in}$ is finite, we find by a compactness argument a labeling of $\\mathbb{Z}^d$ with $\\Sigma_{in}$ such that the event $\\mathcal{F}_r$ that $\\mathcal{A}'$ fails on $\\mathcal{B}(v,r)$ with the restricted labeling has probability bigger than $a$.\nNote that since the labeling is fixed, we have $\\mathcal{F}_r\\supseteq \\mathcal{F}_{r+1}$.\nSet $\\mathcal{F}=\\bigcap_{r\\in \\mathbb{N}} \\mathcal{F}_r$.\nThen by $\\sigma$-additivity we have that the probability of the event $\\mathcal{F}$ is bigger than $a>0$.\nHowever, $F$ produces toast almost surely on $\\mathcal{F}$ and $\\mathcal{A}$ produces $\\Pi$-coloring with probability $1$ for any given toast.\nThis contradicts the fact that probability of $\\mathcal{F}$ is not $0$.\n\n\n\n\\begin{comment}\nIn another words, with probability $1$ any node $v$ can reconstruct bigger and bigger piece of the toast $\\mathcal{D}$ when exploring its neighborhood.\nIn particular, $v$ will eventually became element of some piece of the toast and $\\mathcal{A}$ has to label $v$ with $\\Sigma_{out}$ after finitely many steps.\nThis defines an oblivious local algorithm $\\mathcal{B}$.\n\n\n\nPick a labeling $\\ell\\in [0,1]^{\\mathbb{Z}^d}$ such that $F(\\ell)$ is a $q$-toast and $\\mathcal{A}(\\mathcal{D},\\ell)$ is a $\\Pi$-coloring for $q$-toast $\\mathcal{D}$.\nNote that the picked labeling satisfies these conditions with probability $1$.\n\nDefine $G(\\ell)=\\mathcal{A}(F(\\ell),\\ell)$.\nIt remains to observe that $G$ is an oblivious local algorithm.\nNote that $G(\\ell)({\\bf 0}_d)=\\mathcal{A}(F(\\ell),\\ell)({\\bf 0}_d)$ depends on some $D\\in \\mathcal{D}$ such that ${\\bf 0}_d\\in D$, $\\{E\\subseteq D:E\\in \\mathcal{D}\\}$ and $\\ell\\upharpoonright D$.\nSince $F$ is an oblivious local algorithm, we see that $G$ is too.\nConsequently, $\\Pi$ is in the class $\\mathsf{TAIL}$.\n\\end{comment}\n\nThe idea to prove the converse is to ignore the $q$-toast part and wait for the coding radius.\nSuppose that $\\Pi$ is in the class $\\mathsf{TAIL}$.\nBy the definition there is an oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi$ and an unbounded function $f$ such that\n$$\\P(R_\\mathcal{A}>r)\\le \\frac{1}{f(r)}.$$\nGiven an input labeling of $\\mathbb{Z}^d$ by $\\Sigma_{in}$, we see by the definition of $R_\\mathcal{A}$ that every node $v$ must finish after finitely many steps with probability $1$.\nTherefore it is enough to wait until this finite neighborhood of $v$ becomes a subset of some piece of a given toast $\\mathcal{D}$.\nThis happens with probability $1$ independently of the given toast $\\mathcal{D}$.\n\\end{proof}\n\n\\begin{comment}\nLet $(A,c,\\ell)$ be a triplet as in definition of randomized toast algorithm.\nTake $v\\in A\\setminus \\operatorname{dom}(c)$, and suppose that $v={\\bf 0}_d$.\nOutput $\\mathcal{A}(A,c,\\ell)({\\bf 0}_d)$ if and only if there is $t\\in \\mathbb{N}$ such that $\\mathcal{B}({\\bf 0}_d,t)\\subseteq A$ and $\\ell\\upharpoonright \\mathcal{B}({\\bf 0}_d,t)$ corresponds to some coding neighborhood of $F$.\nThat is to say, any extension $\\ell\\upharpoonright \\mathcal{B}({\\bf 0}_d,t)\\sqsubseteq x\\in [0,1]^{\\mathbb{Z}^d}$ outputs the same color at $F(x)({\\bf 0}_d)$.\nIn that case output $\\mathcal{A}(A,c,\\ell)({\\bf 0}_d)=F(x)({\\bf 0}_d)$.\n\nWe show that $\\mathcal{A}$ produces a $\\Pi$-coloring for $\\mu$-almost every $\\ell\\in [0,1]^{\\mathbb{Z}^d}$ and $1$-toast $\\mathcal{D}$.\nLet $\\ell\\in [0,1]^{\\mathbb{Z}^d}$ be such that $F$ computes a $\\Pi$-coloring an need a finite coding radius for every $v\\in \\mathbb{Z}^d$.\nNote that almost every labeling $\\ell$ satisfies this condition.\nTake a $1$-toast $\\mathcal{D}$.\nIt is easy to see that $\\mathcal{A}(\\mathcal{D},\\ell)$ does not depend on $\\mathcal{D}$ and that $\\mathcal{A}(\\mathcal{D},\\ell)=F(\\ell)$, i.e., $\\mathcal{A}(\\mathcal{D},\\ell)$ is a $\\Pi$-coloring.\nThis finishes the proof.\n\\end{comment}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n \n\n\n\n\n\n\\section{$\\TOAST$: Specific Problems}\\label{sec:SpecificProblems}\n\\todo{explain that some of this is work of others and we show it to substantiate definition of toast}\n\nIn this section we focus on specific problems that are not in the class $\\mathsf{LOCAL}(O(\\log^* n))$ but they admit a solution by a (spontaneous) toast algorithm. Recall, that proving that $\\Pi \\in \\TOAST$ implies that there is an oblivious local algorithm with $1\/r^{1-o(1)}$ tail decay via \\cref{thm:toast_in_tail} and we discuss in \\cref{subsec:descriptive_combinatorics_discussion} that by a result of \\cite{GJKS2} it also implies $\\Pi\\in \\mathsf{BOREL}$. \n\nThe problems we discuss are the following: (a) vertex coloring, (b) edge coloring, (c) tiling the grid with boxes from a prescribed set (this includes, e.g., a perfect matching problem). \n\n\n\n\n\n\nAll these problems were recently studied in the literature from different points of view, see \\cite{ArankaFeriLaci, HolroydSchrammWilson2017FinitaryColoring, GJKS, spencer_personal}.\nIn all cases the main technical tool, sometimes in a slight disguise, is the toast construction.\nOur aim here is to put these problems into a unified context.\nWe note that sometimes the results in the aforementioned literature are stronger or work in a bigger generality (for example, all results discussed here can also be proven to hold on unoriented grids as discussed in \\cite{ArankaFeriLaci}, or the constructed coloring could be aperiodic \\cite{spencer_personal}, etc)\nIn what follows we only give a ``proof by picture'', sketch the description of the corresponding spontaneous toast algorithm and refer the reader to other papers for more details\/stronger results.\n\nThroughout this section we assume that $d>1$.\nWe refer the reader to \\cref{sec:OtherClasses} for definition of classes $\\mathsf{CONTINUOUS}$, $\\mathsf{BOREL}$, etc.\n\n\n\nThe main contribution if this section is the following:\n\\begin{enumerate}\n \\item The fact that the problem of vertex $3$-coloring is in $\\TOAST$ together with the speed-up \\cref{thm:grid_speedup} answers a question (i) from \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n \\item We observe that the problem of edge $2d$-coloring the grid $\\mathbb{Z}^d$ is in $\\TOAST$, thereby answering a question of \\cite{GJKS} coming from the area of descriptive combinatorics. Independently, this was proven in \\cite{ArankaFeriLaci}. \n\\end{enumerate}\n\nIn general, our approach gives that if an LCL $\\Pi$ is in the class $\\TOAST\\setminus \\mathsf{LOCAL}(O(\\log^*(n)))$, then $1\\le {\\bf M}(\\Pi)\\le d-1$ by \\cref{thm:toast_in_tail} and \\cref{thm:grid_speedup}.\nThe examples in Section~\\ref{sec:Inter} show that there are LCLs with ${\\bf M}(\\Pi)=\\ell$ for every natural number $1\\le \\ell\\le d-1$, but for most natural LCLs we do not know the exact value of ${\\bf M}(\\Pi)$, unless $d=2$.\n\n\n\\subsection{Vertex Coloring}\n\\label{subsec:3coloring}\n\nLet $d>1$.\nRecall that $\\Pi^d_{\\chi,3}$ is the \\emph{proper vertex $3$-coloring problem}.\nThis problem was studied extensively in all contexts.\nIt is known that $\\Pi^d_{\\chi,3}\\not\\in \\mathsf{LOCAL}(O(\\log^* n))$ \\cite{brandt_grids}, $\\Pi^d_{\\chi,3}\\not\\in \\mathsf{CONTINUOUS}$ \\cite{GJKS}, $\\Pi^d_{\\chi,3}\\in \\mathsf{BOREL}$ \\cite{GJKS2} and $0<{\\bf M}(\\Pi^d_{\\chi,3})\\le 2$ \\cite{HolroydSchrammWilson2017FinitaryColoring}.\nA version of the following theorem is proven in \\cite{GJKS2}, here we give a simple ``proof by picture''. \n\n\\begin{theorem}[\\cite{GJKS2}]\\label{thm:3in toast}\nLet $d>1$.\nThe LCL $\\Pi^d_{\\chi,3}$ is in the class $\\TOAST$.\n\\end{theorem}\n\n\\begin{proof}[Proof Sketch]\nThe toast algorithm tries to color all nodes in a given piece with two colors, blue and red. However, it can happen that two partial colorings of small pieces that have different parity in a bigger piece have to be connected (see \\cref{fig:3-coloring}). \nIn this case, we take the small pieces colored with incorrect parity and ``cover'' their border with the third color: green. From the inside of a small piece, the green color corresponds to one of the main two colors, while from the outside it corresponds to the other color. This enables us to switch the color parity. \n\\end{proof}\n\nAs a consequence of results from \\cref{sec:Toast} we can now fully determine the moment behavior of the problem which answers a question from \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n\n\\begin{corollary}\nLet $d>1$.\nEvery oblivious local algorithm that solves $\\Pi^d_{\\chi,3}$ has the $1$-moment infinite, but there is an oblivious local algorithm solving $\\Pi^d_{\\chi,3}$ with all $(1-\\varepsilon)$-moments finite. \nIn particular, ${\\bf M}(\\Pi^d_{\\chi,3})=1$. \n\\end{corollary}\n\\begin{proof}\nBy \\cref{thm:3in toast}, $\\Pi^d_{\\chi,3}$ is in the class $\\TOAST$.\nIt follows from \\cref{thm:toast_in_tail} that ${\\bf M}(\\Pi^d_{\\chi,3})\\ge 1$.\nBy the speedup \\cref{thm:grid_speedup} and the fact that $3$-coloring was proven not to be in class $\\mathsf{LOCAL}(O(\\log^* n)) = \\mathsf{TAIL}(\\tow(\\Omega(r))) = \\mathsf{CONTINUOUS}$ \\cite{brandt_grids,HolroydSchrammWilson2017FinitaryColoring,GJKS}, we have that every oblivious local algorithm $\\mathcal{A}$ solving $\\Pi^2_{\\chi,3}$ must have the $1$-moment infinite.\nIn particular, ${\\bf M}(\\Pi^2_{\\chi,3})=1$.\nIt remains to observe that every vertex $3$-coloring of $\\mathbb{Z}^d$ yields a vertex $3$-coloring of $\\mathbb{Z}^2$ by forgetting the last $(d-2)$ dimensions.\nConsequently, every oblivious local algorithm $\\mathcal{A}$ that solves $\\Pi^d_{\\chi,3}$ defines an oblivious local algorithm $\\mathcal{A}'$ that solves $\\Pi^2_{\\chi,3}$ with the same tail behavior.\nThis finishes the proof.\n\\end{proof}\n\nThe $3$-coloring discussion above also implies that the toast itself cannot be constructed faster then with linear tail, otherwise $3$-coloring could be constructed with such a decay, too. \n\\begin{corollary}\n\\label{cor:toast_lower_bound}\nThere is a constant $C > 0$ such that a $q$-toast cannot be constructed in $\\mathsf{TAIL}(C r)$ on any $\\mathbb{Z}^d$.\nHere, by constructing a $q$-toast we mean that the coding radius of each node $u$ is such that $u$ can describe the minimal toast piece $D \\in \\mathcal{D}$ such that $u \\in D$. \n\\end{corollary}\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width = .95\\textwidth]{fig\/3-col-real.eps}\n \\caption{The $\\TOAST$ algorithm is given small pieces of a toast that are already vertex-colored. First step is to match the parity of colors with respect to the bigger piece. Then color the bigger piece.}\n \\label{fig:3-coloring}\n\\end{figure}\n\n\n\\subsection{Edge Coloring}\n\\label{subsec:edgecoloring}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width = .95\\textwidth]{fig\/4-edge-col.png}\n \\caption{The $\\TOAST$ algorithm is given small pieces of a toast that are already edge-colored. First step is to match the parity of colors with respect to the bigger piece. This is done by adding a three layered gadget in the problematic directions. After that the coloring is extended to the bigger piece.}\n \\label{fig:4-edge-coloring}\n\\end{figure}\\todo{dopsat}\nLet $d>1$.\nWe denote as $\\Pi^d_{\\chi',2d}$ the \\emph{proper edge $2d$-coloring problem}.\nObserve that a solution to $\\Pi^d_{\\chi',2d}$ yields a solution to the perfect matching problem $\\Pi^d_{pm}$ (because $\\mathbb{Z}^d$ is $2d$-regular).\nThis shows that $\\Pi^d_{\\chi',2d}$ is not in the class $\\mathsf{LOCAL}(O(\\log^* n))$, because the perfect matching problem clearly cannot be solved with a local algorithm: if it was, we could find a perfect matching in every large enough finite torus, which is clearly impossible since it can have an odd number of vertices \\cite{brandt_grids}. \n\n\n\n\\begin{theorem}[See also \\cite{ArankaFeriLaci}]\\label{thm:edge toast}\nLet $d>1$.\nThen $\\Pi^d_{\\chi',2d}$ is in the class $\\TOAST$.\n\\end{theorem}\n\\begin{proof}[Proof Sketch]\nThe spontaneous toast algorithm is illustrated for $d=2$ in \\cref{fig:4-edge-coloring}.\nThe ``interface'' used by the toast algorithm between smaller and bigger pieces is: smaller pieces pretend they are colored such that in the vertical direction violet and green edges alternate, while in the horizontal direction blue and red edges alternate. If all small pieces in a bigger one have the same parity, this interface is naturally extended to the bigger piece. Hence, we need to work out how to fix pieces with bad horizontal\/vertical parity, similarly to \\cref{thm:3in toast}. \n\n\\cref{fig:4-edge-coloring} shows how to fix the vertical parity, adding the highlighted gadget in the middle picture above and below given pieces keeps its interface the same but switches the parity. The same can be done in the horizontal direction. \n\nTo show the result for $d>2$ it is enough to realize that switching parity in one dimension can be done with a help of one extra additional dimension, while fixing the parity in the remaining dimensions.\nConsequently, this allows to iterate the previous argument at most $d$-many times to fix all the wrong parities.\n\\end{proof}\n\n\n\\begin{corollary}\nLet $d>1$.\nThen $1 \\le {\\bf M}(\\Pi^d_{\\chi',2d}) \\le d-1$. \nIn particular, ${\\bf M}(\\Pi^2_{\\chi',4})=1$.\n\\end{corollary}\n\nWe leave as an open problem to determine the exact value of ${\\bf M}(\\Pi^d_{\\chi',2d})$, where $d>2$.\n\nThe following corollary answers a question of \\cite{GJKS}. See \\cref{sec:OtherClasses} for the definition of class $\\mathsf{BOREL}$. \n\\begin{corollary}\nLet $d\\ge 1$.\nThen $\\Pi^d_{\\chi',2d}$ is in the class $\\mathsf{BOREL}$.\nIn particular, if $\\mathbb{Z}^d$ is a free Borel action and $\\mathcal{G}$ is the Borel graph generated by the standard generators of $\\mathbb{Z}^d$, then $\\chi'_\\mathcal{B}(\\mathcal{G})=2d$, i.e., the Borel edge chromatic number is equal to $2d$.\n\\end{corollary}\n\n\n\\begin{remark}[Decomposing grids into finite cycles]\nBencs, Hru\u0161kov\u00e1, T\u00f3th \\cite{ArankaFeriLaci} showed recently that there is a $\\TOAST$ algorithm that decomposes an unoriented grid into $d$-many families each consisting of finite pairwise disjoint cycles.\nHaving this structure immediately gives the existence of edge coloring with $2d$-colors as well as Schreier decoration.\nSince every $\\TOAST$ algorithms imply Borel solution on graphs that look locally like unoriented grids \\cite{AsymptoticDim}, this result implies that all these problems are in the class $\\mathsf{BOREL}$.\nWe refer the reader to \\cite{ArankaFeriLaci} for more details.\n\\end{remark}\n\n\n\n\n\n\n\n\\begin{comment}\n\\paragraph{Decomposing grids into finite cycles} \\todo{Jsem pro zkratit to na jeden odstavec}\n\nWe sketch how to use toast structure to decompose $\\mathbb{Z}^d$ into $d$-many families each consisting of finite pairwise disjoint cycles, formally, a collection $\\{\\mathcal{C}_i\\}_{i\\in [d]}$ such that $\\mathcal{C}_i$ is family of pairwise vertex disjoint cycles for every $i\\in [d]$ and $\\bigcup_{i\\in [d]} \\mathcal{C}_i=\\mathbb{Z}^d$.\nNote that finding such a decomposition is not an LCL because we do not have bound on the length of the cycles.\nOn the other hand, if we allow infinite paths, then such problem is LCL problem and we denote it as $\\Pi^d_{\\operatorname{cyc},d}$.\nMoreover, it is easy to see that since we work with oriented grids $\\Pi^d_{\\operatorname{cyc},d}$ is in the class $\\mathsf{LOCAL}(O(1))$.\n\nWe show that there is a spontaneous extension function $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is always a $\\Pi^d_{\\operatorname{cyc},d}$-coloring with the additional property that all the cycles in $\\mathcal{A}(\\mathcal{D})$ are finite for every $37$-toast $\\mathcal{D}$.\nMoreover, it is easy to see that our toastable algorithm also works for non-oriented grids, where the LCL problem $\\Pi^d_{\\operatorname{cyc},d}$ is not trivial.\nThe combinatorial construction is inspired by result of Bencs, Hruskova and Toth \\cite{ArankaFeriLaci} who originally proved similar result for $d=2$.\nAfter reading their paper we realized that the argument goes through for any dimension $d>1$ and we refer the reader to their paper for a detailed proof of this fact.\nHere we only sketch the main idea.\n\n\\begin{theorem}\nLet $d>1$.\nThen there is a spontaneous extension function $\\mathcal{A}$ such that $\\mathcal{A}(\\mathcal{D})$ is always a $\\Pi^d_{\\operatorname{cyc},d}$-coloring for every $37$-toast $\\mathcal{D}$ with the additional property that all the cycles in $\\mathcal{A}(\\mathcal{D})$ are finite.\n\\end{theorem}\n\n\\begin{remark}\nA straightforward use of the toastable algorithm together with the fact that Borel graphs with finite Borel asymptotic dimension admit Borel toast structure (see~\\cite{AsymptoticDim}) gives a decomposition of any Borel graph that looks locally like a $d$-dimensional grid (oriented or non-oriented) into $d$-many Borel families of edges each consisting of pairwise disjoint cycles.\nSince each such cycle is necessarily of even length, we get immediately that such graphs admit Borel Schreier decoration, Borel balanced orientation and their Borel chromatic index is $2d$.\n\\end{remark}\n\n\\begin{proof}[Proof Sketch]\nLet $d>1$.\nIt is easy to see that decomposing $\\mathbb{Z}^d$ into $d$ families of cycles (possibly infinite) is the same as coloring the edges with $d$-colors, such that every vertex $v\\in \\mathbb{Z}^d$ has exactly two adjacent edges of each color.\nLet $\\{\\alpha_1,\\dots,\\alpha_d\\}$ be the set of the colors.\n\nAn \\emph{(edge) coloring} of a connected finite set $A\\subseteq \\mathbb{Z}^d$ is a coloring of all edges that are adjacent to $A$.\nWe say that an edge $e$ is a \\emph{boundary edge}, if only one vertex of $e$ is contained in $A$.\nA coloring $c$ is \\emph{suitable}, if after choosing a reference point and orientation of $A$ we have that parallel boundary edges in the $i$-th direction are assigned the color $\\alpha_i$.\n\nThere is a constant $q\\in \\mathbb{N}$ such that for every $i\\not=j\\in [d]$, finite connected set $A\\subseteq \\mathbb{Z}^d$ and a suitable edge coloring $c$, there is a suitable edge coloring $c'$ of $A':=\\mathcal{B}(A,q)$ that extends $c$ such that \n\\begin{itemize}\n \\item every $c'$-cycle of color $\\alpha_i$ that contains a vertex from $A$ is finite for every $i\\in [d]$,\n \\item the $i$-th and $j$-th directions of $A$ and $A'$ are swapped, e.g., parallel $A'$-boundary edges in the $i$-th direction from the point of view of $A$ are $c'$-colored with $\\alpha_j$,\n \\item other directions remain the same.\n\\end{itemize}\nIt is easy to see that iterating this extension possibly $2d$-many times produces a spontaneous toast algorithm $\\mathcal{A}$ with all the desired properties.\n\nIt remains to sketch how to find such extension $c'$.\nFirst, since $c$ is suitable and we want to swap $i$-th and $j$-th direction, we may assume that $d=2$.\nSecond, we may assume that $A$ is a union of rectangles, such that each side has even number of vertices adjacent to a perpendicular boundary edge.\nThis could be done as follows, pick a reference point and consider decomposition of $\\operatorname{Cay}(\\mathbb{Z}^2)$ into a square of a fixed even number of vertices on a side.\nTake the union of all such squares that intersect $A$.\nThird, see Figure~\\ref{fig:SwapingColors} how to swap colors and close all cycles coming out of $A$.\nThis finishes the sketch.\n\\end{proof}\n\\end{comment}\n\n\n\n\n\\subsection{Tiling with Rectangles}\n\\label{subsec:rectangle_tiling}\n\n\nLet $d>1$ and $\\mathcal{R}$ be a finite collection of $d$-dimensional rectangles.\nWe denote as $\\Pi^d_{\\mathcal{R}}$ the problem of \\emph{tiling $\\mathbb{Z}^d$ with elements of $\\mathcal{R}$}.\nIt is not fully understood when the LCL $\\Pi^d_{\\mathcal{R}}$ is in the class $\\mathsf{LOCAL}(O(\\log^* n))$ (or equivalently $\\mathsf{CONTINUOUS}$).\nCurrently the only known condition implying that $\\Pi^d_{\\mathcal{R}}$ is \\emph{not} in $\\mathsf{LOCAL}(O(\\log^* n))$ is when GCD of the volumes of elements of $\\mathcal{R}$ is not $1$.\nMinimal unknown example is when $\\mathcal{R}=\\{2\\times 2, 2\\times 3,3\\times 3\\}$. \nFor the other complexity classes the situation is different.\nGiven a collection $\\mathcal{R}$ we define a condition \\eqref{eq:GCD} as\n\\begin{equation}\n\\label{eq:GCD}\\tag{$\\dagger$}\n \\text{GCD of side lengths in the direction of} \\ i \\text{-th generator is} \\ 1 \\ \\text{for every} \\ i\\in [d].\n\\end{equation}\nIt follows from the same argument as for vertex $2$-coloring, that if $\\mathcal{R}$ does not satisfy \\eqref{eq:GCD}, then $\\Pi^d_{\\mathcal{R}}$ is a global problem, i.e., it is not in any class from \\cref{fig:big_picture_grids}.\nOn the other hand, the condition \\eqref{eq:GCD} is sufficient for a spontaneous toast algorithm. For example, this implies $\\Pi^2_{\\{2\\times 3,3\\times 2\\}}\\in \\mathsf{BOREL} \\setminus \\mathsf{CONTINUOUS}$.\nA version of the following theorem was announced in \\cite{spencer_personal}, here we give a simple ``proof by picture''.\n\n\\begin{theorem}[\\cite{spencer_personal, stephen}]\n\\label{thm:tilings}\nLet $d>1$ and $\\mathcal{R}$ be a finite collection of $d$-dimensional rectangles that satisfy \\cref{eq:GCD}.\nThen $\\Pi^d_{\\mathcal{R}}$ is in the class $\\TOAST$.\nConsequently, ${\\bf M}(\\Pi^d_{\\mathcal{R}})\\ge 1$.\n\\end{theorem}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=.9\\textwidth]{fig\/rectangle_filling.eps}\n \\caption{The $\\TOAST$ algorithm is given small pieces of a toast that are already tiled. First step is to match the parity of the tiling with respect to the bigger piece. Then tile the bigger piece.}\n \\label{fig:Tiling}\n\\end{figure}\n\n\\begin{proof}[Proof Sketch]\nThe two-dimensional version of the construction is shown in \\cref{fig:Tiling}. The interface of pieces this time is that they cover themselves with boxes of side length $L$ which is the least common multiple of all the lengths of basic boxes. As in the previous constructions, we just need to work out how to ``shift'' each piece so that all small pieces agree on the same tiling with boxes of side length $L$ which can then be extended to the whole piece. This can be done as in \\cref{fig:Tiling}. \n\\end{proof}\n\n\n\\begin{corollary}\nLet $d=2$ and $\\mathcal{R}$ be a finite collection of $2$-dimensional rectangles.\nThen $\\Pi^2_{\\mathcal{R}}$ is either in the class $\\mathsf{LOCAL}(O(\\log^* n))$, satisfies ${\\bf M}(\\Pi^2_{\\mathcal{R}})=1$, or is global (that is, in no class in \\cref{fig:big_picture_grids}).\n\\end{corollary}\n\nFor any $d$, \\cref{thm:tilings} implies that if a tiling problem is ``hard'' (in no class in \\cref{fig:big_picture_grids}), it is because of a simple reason: a projection to one dimension gives a periodic problem. It would be interesting to understand whether there are problems that are hard but their solutions do not have any periodicity in them (see \\cref{sec:open_problems}). \n\nWe do not know if there is $\\mathcal{R}$ that would satisfy $1< {\\bf M}(\\Pi^d_{\\mathcal{R}})\\le d-1$.\nLet $d\\ge \\ell>0$ and denote as $\\mathcal{R}_\\ell$ the restriction of rectangles to first $\\ell$-dimensions.\nAs in the case of proper $3$-vertex coloring problem, we have that ${\\bf M}(\\Pi^d_{\\mathcal{R}})\\le {\\bf M}(\\Pi^\\ell_{\\mathcal{R}_{\\ell}})$.\nIn particular, if, e.g., $\\Pi^d_{\\mathcal{R}_2}$ is not in the class $\\mathsf{LOCAL}(O(\\log^* n))$, then ${\\bf M}(\\Pi^d_{\\mathcal{R}})\\le 1$.\n\n\\begin{remark}\nIn the context of non-oriented grids we need to assume that $\\mathcal{R}$ is invariant under swapping coordinates. \nSimilar argument, then shows that $\\Pi^d_\\mathcal{R}\\in \\TOAST$ if and only if $\\mathcal{R}$ satisfies \\eqref{eq:GCD}.\nCombined with the fact that there is a Borel toast on such Borel graphs (see~\\cite{AsymptoticDim}), we deduce that \\eqref{eq:GCD} for $\\mathcal{R}$ is equivalent with existence of Borel tiling with rectangles from $\\mathcal{R}$.\n\\end{remark}\n\n\n\\todo{2x3 a 3x2 se daji poskladat z elek, coz odpovida nejake partition kruznice}\n\n\n\n\n\n\\paragraph{Perfect matching}\n\nLet $d>1$.\nA special case of a tiling problem is the \\emph{perfect matching problem} $\\Pi^d_{pm}$.\nSince there are arbitrarily large tori of odd size we conclude that $\\Pi^d_{pm}\\not \\in \\mathsf{LOCAL}(O(\\log^* n))$.\nConsequently, $\\Pi^d_{pm}\\not \\in \\mathsf{TAIL}(\\tow(\\Omega(r)))$ by \\cref{thm:tail=uniform}, and it follows from the Twelve Tile Theorem of Gao et al \\cite[Theorem 5.5]{GJKS} that $\\Pi^d_{pm}\\not \\in \\mathsf{CONTINUOUS}$.\nSimilarly as with the vertex $3$-coloring problem, the proof of $\\Pi^d_{pm}\\in \\mathsf{BOREL}$, where $d>1$, from \\cite{GJKS2} uses the toast construction and, in our notation, yields the following result. \n\n\\begin{theorem}[\\cite{GJKS2}]\\label{thm:perfect_matching_intoast}\nLet $d>1$.\nThe LCL $\\Pi^d_{pm}$ is in the class $\\TOAST$.\n\\end{theorem}\n\n\\begin{corollary}\nLet $d>1$.\nThen ${\\bf M}(\\Pi^d_{pm})\\ge 1$ and every ffiid $F$ that solves $\\Pi^d_{\\chi,3}$ has the $(d-1)$-moment infinite.\nIn particular, ${\\bf M}(\\Pi^2_{pm})=1$.\n\\end{corollary}\n\nWe leave as an open problem to determine the exact value of ${\\bf M}(\\Pi^d_{pm})$, where $d>2$.\nIn fact, this problem is similar to intermediate problems that we study in \\cref{sec:Inter}. \n\n\n\n\n\n\n\n\n\n\n\n\\renewcommand{\\pm}[2]{\\operatorname{PM}^{#1, #2}_\\boxtimes}\n\\renewcommand{\\pm}[1]{\\operatorname{PM}^{#1}_\\boxtimes}\n\\newcommand{\\mathbb{Z}^d_\\boxtimes}{\\mathbb{Z}^d_\\boxtimes}\n\n\\section{Intermediate Problems}\\label{sec:Inter}\n\nIn this section we construct local problems with $\\mathsf{TAIL}$ complexities $\\Theta(r^{j-o(1)})$ for $1 \\le j \\le d-1$. \nWe first define the problem $\\pm{d}$: this is just a perfect matching problem, but in the graph $\\mathbb{Z}^d_\\boxtimes$ where two different nodes $u = (u_1, \\dots, u_d)$ and $v = (v_1, \\dots, v_d)$ are connected with an edge if $\\sum_{1 \\le i \\le d} |u_i - v_i| \\le 2$. That is, we also add some ``diagonal edges'' to the grid. We do not know whether perfect matching on the original graph (Cayley graph of $\\mathbb{Z}^d$) is also an intermediate problem (that is, not $\\mathsf{LOCAL}$ but solvable faster than $\\TOAST$). \n\nWe will show that the tail complexity of $\\pm{d}$ is $\\Theta(r^{d - 1 - o(1)})$. \nAfter this is shown, the hierarchy of intermediate problems is constructed as follows. \nConsider the problem $\\pm{d,\\ell}$. This is a problem on $\\mathbb{Z}^d$ where one is asked to solve $\\pm{d}$ in every hyperplane $\\mathbb{Z}^\\ell$ given by the first $\\ell$ coordinates in the input graph. \nThe $\\mathsf{TAIL}$ complexity of $\\pm{d,\\ell}$ is clearly the same as the complexity of $\\pm{\\ell}$, that is, $\\Theta(r^{\\ell-1-o(1)})$. \nHence, we can focus just on the problem $\\pm{d}$ from now on. \n\n\\begin{claim}\\label{cl:PM is global}\nLet $2\\le d\\in \\mathbb{N}$.\nThen the local coloring problem $\\pm{d}$ is not in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\nVia \\cref{thm:grid_speedup} this is equivalent to $\\pm{d} \\not\\in \\mathsf{TAIL}(Cr^{d-1})$ for large enough $C$. \n\\end{claim}\n\nThis follows from the fact that if there is a local algorithm, then in particular there exists a solution on every large enough torus. But this is not the case if the torus has an odd number of vertices. \n\nSecond, we need to show how to solve the problem $\\pm{d}$ faster than with a $\\TOAST$ construction algorithm. \n\n\\begin{restatable}{theorem}{intermediate}\n\\label{thm:intermediate}\nThe problem $\\pm{d}$ can be solved with an oblivious local algorithm of tail complexity $r^{d-1-O(\\frac{1}{\\sqrt{\\log r}})}$. \n\\end{restatable}\n\nHere we give only a proof sketch, the formal proof is deferred to \\cref{app:intermediate_problems}. \n\n\\begin{proof}[Proof Sketch]\nWe only discuss the case $d=3$ (and the actual proof differs a bit from this exposition due to technical reasons). As usual, start by tiling the grid with boxes of side length roughly $r_0$. \nThe \\emph{skeleton} $\\mathcal{C}_0$ is the set of nodes we get as the union of the $12$ edges of each box. \nIts \\emph{corner points} $X_0 \\subseteq \\mathcal{C}_0$ is the union of the $8$ corners of each box. \n\n\nWe start by matching vertices inside each box that are not in the skeleton $\\mathcal{C}_0$. It is easy to see that we can match all inside vertices, maybe except of one node per box that we call a \\emph{pivot}. We can ensure that any pivot is a neighbor of a vertex from $\\mathcal{C}_0$. \n\nWe then proceed similarly to the toast construction of \\cref{thm:toast_in_tail} and build a sequence of sparser and sparser independent sets that induce sparser and sparser subskeletons of the original skeleton.\nMore concretely, we choose a sequence $r_k = 2^{O(k^2)}$ and in the $k$-th step we construct a new set of corner points $X_k \\subseteq X_{k-1}$ such that any two corner points are at least $r_k$ hops apart but any point of $\\mathbb{Z}^d$ has a corner point at distance at most $2r_k$ (in the metric induced by $\\mathbb{Z}^d$). \nThis set of corner points $X_k$ induces a subskeleton $\\mathcal{C}_{k} \\subseteq \\mathcal{C}_{k-1}$ by connecting nearby points of $X_k$ with shortest paths between them in $\\mathcal{C}_{k-1}$. \n\nAfter constructing $\\mathcal{C}_k$, the algorithm solves the matching problem on all nodes of $\\mathcal{C}_{k-1} \\setminus \\mathcal{C}_k$. This is done by orienting those nodes towards the closest $\\mathcal{C}_k$ node in the metric induced by the subgraph $\\mathcal{C}_k$. \nThen, the matching is found on connected components that have radius bounded by $r_k$. There is perhaps only one point in each component, a pivot, that neighbors with $\\mathcal{C}_{k}$ and remains unmatched. \n\nAfter the $k$-th step of this algorithm, only nodes in $\\mathcal{C}_k$ (and few pivots) remain unmatched. Hence, it remains to argue that only roughly a $1\/r_k^{d-1}$ fraction of nodes survives to $\\mathcal{C}_k$. This should be very intuitive -- think about the skeleton $\\mathcal{C}_k$ as roughly looking as the skeleton generated by boxes of radius $r_k$. This is in fact imprecise, because in each step the distance between corner points of $X_k$ in $\\mathcal{C}_k$ may be stretched by a constant factor with respect to their distance in the original skeleton $\\mathcal{C}_{k-1}$. This is the reason why the final volume of $\\mathcal{C}_k$ is bounded only by $2^{O(k)} \/ r_k^{d-1}$ and by a similar computation as in \\cref{thm:toast_in_tail} we compute that the tail decay of the construction is $2^{O(\\sqrt{\\log r})}\/{r^{d-1}}$. \n\\end{proof}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Other Classes of Problems}\\label{sec:OtherClasses}\n\nIn this mostly expository section we discuss other classes of problems from \\cref{fig:big_picture_grids} and how they relate to our results from previous sections. \n\n\n\n\\subsection{$\\mathsf{CONTINUOUS}$}\\label{subsec:CONT}\n\nHere we explain the fact that an LCL $\\Pi$ on $\\mathbb{Z}^d$ is in the class $\\mathsf{LOCAL}(O(\\log^*(n)))$ if and only if it is in the class $\\mathsf{CONTINUOUS}$.\nFor $\\mathbb{Z}^d$, the model $\\mathsf{CONTINUOUS}$ was extensively studied by Gao, Jackson, Krohne, Tucker-Drob, Seward and others, see \\cite{GaoJackson,GJKS,ST,ColorinSeward}.\nIn fact, the connection with local algorithms is implicit in \\cite{GJKS} under the name Twelve Tiles Theorem \\cite[Theorem~5.5]{GJKS}.\nRecently, independently Bernshteyn \\cite{Bernshteyn2021local=cont} and Seward \\cite{Seward_personal} showed that the same holds in bigger generality for all Cayley graphs of finitely generated groups.\nOur aim here is to formally introduce the model $\\mathsf{CONTINUOUS}$ and give an informal overview of the argument for $\\mathbb{Z}^d$ that follows the lines of Gao et al \\cite{GJKS}.\nThe reasons not to simply put reference to the aforementioned papers are two.\nFirst, we work with slightly more general LCLs that have inputs.\nSecond, our goal is to introduce all classes in a way accessible for a broader audience. \n\nThe setup of the $\\mathsf{CONTINUOUS}$ model is somewhat similar to the deterministic $\\mathsf{LOCAL}$ model: there is no randomness in the definition and there is a certain symmetry breaking input on the nodes. Specifically, each vertex is labelled with one label from the set $\\{0, 1\\}$; the space of all such labelings is denoted as $2^{\\mathbb{Z}^d}$. \nHowever, for the same reason as why we insist on the uniqueness of the identifiers in the $\\mathsf{LOCAL}$ model, we insist that the input $\\{0,1\\}$ labeling breaks all symmetries. Namely, we require that there is no $g$ such that shifting all elements by $g$ preserves the labeling. We write $\\operatorname{Free}(2^{\\mathbb{Z}^d}) \\subseteq 2^{\\mathbb{Z}^d}$ for this subspace. \nA \\emph{continuous algorithm} is a pair $\\mathcal{A}=(\\mathcal{C},f)$, where $\\mathcal{C}$ is a collection of $\\{0,1\\} \\times \\Sigma_{in}$-labeled neighborhoods and $f$ is a function that assigns to each element of $\\mathcal{C}$ an output label. It needs to satisfy the following. Suppose we fix a node $v$ and the input $\\{0,1\\}$-labeling $\\Lambda$ of $\\mathbb{Z}^d$ satisfies $\\Lambda \\in \\operatorname{Free}(2^{\\mathbb{Z}^d})$. If we let $v$ explore its neighborhood with an oblivious local algorithm, then we encounter an element from $\\mathcal{C}$ after finitely many steps. \nWe write $\\mathcal{A}(\\Lambda)(v)$, or simply $\\mathcal{A}(v)$, when $\\Lambda$ is understood, for the corresponding output.\n\n\\begin{comment}\n\\vspace*{5cm}\nTo break symmetries on $\\mathbb{Z}^d$ in the setting of random processes one considers iid labels on vertices.\nThe event that some symmetries remain has measure $0$ and we simply ignore such labelings.\nIn the continuous settings we label vertices with $2:=\\{0,1\\}$.\nThe space of all such labelings is denoted as $2^{\\mathbb{Z}^d}$ and we write $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ for the space of labelings that break all symmetries.\nA \\emph{continuous algorithm} is a pair $\\mathcal{A}=(\\mathcal{C},g)$, where $\\mathcal{C}$ is a collection of ${0,1},\\Sigma_{in}$-labeled neighborhoods and $g$ is a function that assigns to each element of $\\mathcal{C}$ an output label, such that if we let a node $v$ explore its labeled neighborhood $x$, then it encounters element from $\\mathcal{C}$ after finitely many steps whenever $x\\in \\Sigma_{in}^{\\mathbb{Z}^d}\\times \\operatorname{Free}(2^{\\mathbb{Z}^d})$, i.e., $x$ breaks all symmetries.\nWe write $\\mathcal{A}(x)(v)$ for the corresponding output.\n\\end{comment}\n\n\\begin{remark}\nThis is equivalent with the definition given in e.g. \\cite{GJKS} in the same way as the correspondence between oblivious local algorithms and ffiid, see \\cref{subsec:subshifts}.\nThe fact that the collection $\\mathcal{C}$ is dense, i.e., every node finishes exploring its neighborhoods in finite time, means that $\\mathcal{A}({\\bf 0})$ is \\emph{continuous} with respect to the product topology and defined in the whole space $\\operatorname{Free}(2^{\\mathbb{Z}^d})$.\nMoreover, the value computed at $v$ does not depend on the position of $v\\in \\mathbb{Z}^d$ but only on the isomorphism type of the labels, this is the same as saying that $\\mathcal{A}(\\_)$ is \\emph{equivariant}.\n\\end{remark}\n\nWe say that an LCL $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ if there is a continuous algorithm $\\mathcal{A}$ that outputs a solution for every $\\Lambda \\in \\operatorname{Free}(2^{\\mathbb{Z}^d})$.\nNote that, in general, we cannot hope for the diameters of the neighborhoods in $\\mathcal{C}$ to be uniformly bounded.\nIn fact, if that happens then $\\Pi \\in \\mathsf{LOCAL}(O(1))$. \nThis is because there are elements in $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ with arbitrarily large regions with constant $0$ label.\nApplying $\\mathcal{A}$ on that regions gives an output label that depends only on $\\Sigma_{in}$ after bounded number of steps, hence we get a local algorithm with constant local complexity.\n\nThe main idea to connect this context with distributed computing is to search for a labeling in $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ where any $\\mathcal{A}$ necessarily has a uniformly bounded diameter.\nSuch a labeling does exist!\nIn \\cite{GJKS} these labelings are called \\emph{hyperaperiodic}.\nTo show that hyperaperiodic labelings exist on $\\mathbb{Z}^d$ is not that difficult, unlike in the general case for all countable groups \\cite{ColorinSeward}. \nIn fact there is enough flexibility to encode in a given hyperaperiodic labeling some structure in such a way that the new labeling remains hyperaperiodic.\n\nWe now give a definition of a hyperaperiodic labeling. \nLet $\\Lambda \\in 2^{\\mathbb{Z}^d}$ be a labeling of $\\mathbb{Z}^d$.\nWe say that $\\Lambda$ is \\emph{hyperaperiodic} if for every $g\\in \\mathbb{Z}^d$ there is a finite set $S\\subseteq \\mathbb{Z}^d$ such that \nfor every $t\\in \\mathbb{Z}^d$ we have \n$$\\Lambda(t+S) \\not= \\Lambda(g+t+S)$$\nThis notation means that the $\\Lambda$-labeling of $S$ differs from the $\\Lambda$-labeling of $S$ shifted by $g$ everywhere. \nIn another words, given some translation $g\\in \\mathbb{Z}^d$, there is a uniform witness $S$ for the fact that $\\Lambda$ is nowhere $g$-periodic.\nIt follows from a standard compactness argument that any continuous algorithm $\\mathcal{A}$ used on a hyperaperiodic labeling finishes after number of steps that is uniformly bounded by some constant that depends only on $\\mathcal{A}$. \n\n\n\n\\begin{comment}\n\\vspace*{5cm}\n\nLet $\\Lambda \\in 2^{\\mathbb{Z}^d}$ be a labeling of $\\mathbb{Z}^d$.\nWe say that $\\Lambda$ is \\emph{hyperaperiodic} if for every $g\\in \\mathbb{Z}$ there is a finite set $S\\subseteq \\mathbb{Z}^d$ such that \nfor every $s\\in \\mathbb{Z}^d$ we have \n$$x(s+t)\\not=x(g+s+t)$$\nfor some $t\\in S$.\nIn another words, given a shift $g\\in \\mathbb{Z}^d$ there is a uniform witness $S$ for the fact that $x$ is nowhere $g$-periodic.\nIt follows from some standard compactness argument that every hyperaperiodic element needs uniformly bounded neighborhoods for each continuous algorithm $\\mathcal{A}$, the bound of course depends on the algorithm. \n\\end{comment}\n\n\\begin{theorem}[\\cite{GJKS,Bernshteyn2021local=cont}]\nLet $d>0$.\nThen an LCL $\\Pi$ on $\\mathbb{Z}^d$ is in the class $\\mathsf{CONTINUOUS}$ if and only if it is in the class $\\mathsf{LOCAL}(O(\\log^*(n)))$.\n\\end{theorem}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=.9\\textwidth]{fig\/continuous.eps}\n \\caption{On the left side we have examples of a labeled local tiling configurations compared with a possible continuity window $w$. These are then sparsely embedded into $\\mathbb{Z}^2$.}\n \\label{fig:my_label}\n\\end{figure}\n\n\\begin{proof}[Proof Sketch]\nOne direction was observed by Bernshteyn \\cite{Bernshteyn2021local=cont} and follows from the fact that there is continuous algorithm that produces a maximal $r$-independent set for every $r\\in \\mathbb{N}$ (recall \\cref{thm:classification_of_local_problems_on_grids}).\nThis is implicit, even though in a different language, already in Gao and Jackson \\cite{GaoJackson}.\n\nNext, suppose that $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ and write $\\mathcal{A}$ for the continuous algorithm for $\\Pi$. \nWe will carefully construct a certain hyperaperiodic element $\\Lambda$ and let $\\mathcal{A}$ solve $\\Pi$ on it. \nWe know that there is a certain number $w$ such that $\\mathcal{A}$ solves $\\Pi$ on $\\Lambda$ with local complexity $w$. \nWe note that even though we do not know the value $w$ in advance, it depends on $\\Lambda$ that is due to be constructed, the flexibility of the construction guarantees to ``encode'' in $\\Lambda$ all possible scenarios.\nSee the discussion at the end of the argument. \n\nThe overall approach is as follows. \nOur local algorithm $\\mathcal{A}'$ first tiles the plane with tiles of length much bigger than $w$ with $O(\\log^* n)$ local complexity. We cannot tile the plane in a periodic manner in $O(\\log^* n)$ steps, but it can be done with a constant number of different tiles such that the side length of each tile is either $r$ or $r+1$ for some $r \\gg w$ (this construction is from \\cite[Theorem~3.1]{GaoJackson}, see also the proof of \\cref{thm:grid_speedup}). \nFor each tile $T$, $\\mathcal{A}'$ labels it with certain $\\{0,1\\}$ labeling. \nThis labeling depends only on the shape of the tile.\nWe make sure that not only $T$ with its labeling, but also $T$ together with all possible local configurations of tiles around it and corresponding labelings appears in $\\Lambda$.\nIn particular, as $r\\gg t,w$ ($t$ is the radius necessary to check the solution $\\Pi$), we have that the $(w+t)$-hop neighborhood of any node is labelled in a way that also appears somewhere in the labeling $\\Lambda$.\nAfter labeling each tile, $\\mathcal{A}'$ simulates $\\mathcal{A}$. We know that $\\mathcal{A}$ finishes after $w$ steps on $\\Lambda$ and is correct there, hence $\\mathcal{A}'$ is correct and has local complexity $O(\\log^* n)$. What remains is to describe is how to get the labeling of each tile. \n\nWe use the fact that there is a lot of flexibility when constructing a hyperaperiodic labeling. \nWe define the hyperaperiodic labeling $\\Lambda$ of the input graph $\\mathbb{Z}^d$ in such a way that, after fixing $r$, each one of constantly many local configurations of $\\{r,r+1\\}$-tiles are embedded in $\\Lambda$, each one in fact in many copies that exhaust possible $\\Sigma_{in}$-labelings.\n\nThis is done as follows.\nFix any hyperaperiodic element $\\Lambda'$, these exist by \\cite[Lemma~2.8]{GJKS}.\nGiven a tile $T$, that depends on $r$, $\\{0,1\\}$-label it by any $T$-pattern that appears in $\\Lambda'$.\nConsider the collection of all possible local configurations of such labeled tiles that are moreover labeled by $\\Sigma_{in}$.\nThis collection is finite. Since we do not know $w$ in advance, do this for any value of $r$. This yields a countable collection of possible configurations. \nEmbed this collection sparsely into $\\Lambda'$, \\cref{fig:my_label}. Extend the $\\Sigma_{in}$-labeling arbitrarily. \nThis defines $\\Lambda$.\nIt is routine to check that $\\Lambda$ is hyperaperiodic, see \\cite[Lemma~5.9]{GJKS}.\n\n\n\nIt remains to recall that since $\\Lambda$ is hyperaperiodic, there is $w\\in \\mathbb{N}$ such that $\\mathcal{A}$ has uniformly bounded locality by $w$ on $\\Lambda$.\nPick $r\\gg w,t$.\nThis works as required since any $w+t$-hop neighborhood $\\mathcal{A}'$ can encounter is present in $\\Lambda$.\n\\end{proof}\n\n\\begin{comment}\n\\vspace*{4cm}\nNext, suppose that $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ and write $\\mathcal{A}$ for the continuous algorithm.\nA high-level strategy for $d=2$ is as follows.\nThere are 12 type of tiles that are parametrized by natural numbers $n,p,q\\in \\mathbb{N}$, i.e., for fixed $n,p,q$ there are $12$ tiles.\nIf $n,p,q$ are fixed, then there is a local deterministic algorithm that tiles given, large enough, torus with these tiles in such a way that the boundary parts matches.\\todo{picture here ad reference to the nice pictures in \\cite{GJKS}}\nIt is not clear at this point how large parameters we should choose.\nThe trick is to encode all such possible tiles for all $n,p,q\\in \\mathbb{N}$ in one hyperaperiodic element $x$.\nThe encoding of the tiles in $x$ guarantees that whenever the regions \\todo{picture} are the same for the fixed parameters $n,p,q$, then $\\mathcal{A}$ outputs the same label.\nSince $x$ is hyperaperiodic, we have that $\\mathcal{A}$ stops after $T$-many steps to output $\\mathcal{A}(x)(v)$ for some $T\\in \\mathbb{N}$.\nWe choose $n,p,q$ according to $T$ in such a way that for any vertex at any region there is a tile such that the $T$-neighborhood of this vertex is fully contained in the tile.\nConsequently, this defines local rule once a given, large enough, torus is tiled into these tiles of 12 types.\n\nThis description exactly follows the proof of (1) $\\rightarrow$ (2) in \\cite[Theorem~5.5]{GJKS}, where (1) is the same as saying that there is a continuous algorithm and (2) fixes the coloring of the tiles for some parameters $n,p,q$.\nHowever, we note that we are skipping some technical details about additional parameter $w$ etc.\nThe additional step to produce a local algorithm by tiling a large enough torus into these 12 tiles is standard.\nIn fact, this is proved in \\cite[Theorem~5.10]{GJKS}, again in different language, to show (2) $\\rightarrow$ (1) in \\cite[Theorem~5.5]{GJKS}.\n\nThe main difference that we encounter are the input labels.\nWe explain how to handle them and hint, simultaneously, what we mean by encoding tiles in hyperaperiodic element, again ignoring a technical parameter $w$, see \\cite[Figures~13,14,etc]{GJKS}.\nFix a hyperaperiodic element $y$.\nFor $n,p,q\\in \\mathbb{N}$ and a tile $G$, see \\eqref{fig:???}, assign to each region labels from $\\{0,1\\}$ in such a way that (a) same regions get the same pattern, (b) this pattern is part of $y$.\nConsider moreover any possible input labeling of such $\\{0,1\\}$-labeled tiles.\nNow place these labeled tiles, for every $n,p,q\\in \\mathbb{N}$, sparsely in $y$ (Figure ???) and extend the input labeling on them to the whole $\\mathbb{Z}^2$, this defines $x$.\nIt is not hard to see that $x$ remains hyperaperiodic.\n\nIn higher dimension, i.e., when $d>2$, the same argument works as well with the only difference that there are more than 12 types of tiles.\n\\end{comment}\n\n\n\n\n\n\n\n\n\n\\subsection{Descriptive Combinatorics Classes}\n\\label{subsec:descriptive_combinatorics_discussion}\n\nIn this section we define the classes of problems that are studied in descriptive combinatorics.\nThere are two equivalent ways how the notion of a \\emph{Borel graph} can be defined.\nThe less abstract one is as follows.\nConsider the unit interval $[0,1]$ with the $\\sigma$-algebra of Borel sets, that is the smallest algebra of subsets that contains sub-intervals and is closed under countable unions and complements.\nLet $A,B\\subseteq [0,1]$ be Borel sets and consider a bijection $\\varphi:A \\to B$ between them.\nWe say that $\\varphi$ is \\emph{Borel measurable} if $\\varphi^{-1}(D\\cap B)$ is a Borel set for every Borel set $D$.\nNow, let $\\mathcal{G}$ be a graph with vertex set $[0,1]$ and assume that the degree of every vertex is bounded by some $\\Delta<\\infty$.\nWe say that $\\mathcal{G}$ is a \\emph{Borel graph} if the edge relation is induced by Borel bijections, that is there are sets $A_1,\\dots, A_\\ell$, $B_1,\\dots, B_\\ell$ and Borel bijections $\\varphi_1,\\dots, \\varphi_\\ell$ such that $(x,y)$ is an edge in $\\mathcal{G}$ if and only if there is $1\\le i\\le \\ell$ such that either $\\varphi_i(x)=y$ or $\\varphi_i(y)=x$.\n\nThe more abstract way how to define Borel graphs uses the notion of a \\emph{standard Borel space}.\nFor example $[0,1]$ with the $\\sigma$-algebra of Borel sets, when we forget that it comes from the standard topology, is a standard Borel space.\nFormally, a standard Borel space is a pair $(X,\\mathcal{B})$, where $X$ is a set and $\\mathcal{B}$ is a $\\sigma$-algebra that comes from some separable and complete metric structure on $X$.\nA uniformly bounded degree graph $\\mathcal{G}$ on $(X,\\mathcal{B})$ is Borel if it satisfies the same condition as in the case of the particular space $[0,1]$.\n\nIt is easy to see that with this definition all the basic constructions from graph theory are \\emph{relatively} Borel, e.g. if $A$ is a Borel set of vertices, then $N_\\mathcal{G}(A)$ the neighbors of elements in $A$ form a Borel set.\nAlso note that since we assume that the degree is uniformly bounded, every connected component in $\\mathcal{G}$ is at most countable.\nTypical questions that are studied in this context are Borel versions of standard notions, e.g., \\emph{Borel} chromatic number, \\emph{Borel} chromatic index, \\emph{Borel} perfect matching etc.\nWe refer the reader to \\cite{pikhurko2021descriptive_comb_survey,kechris_marks2016descriptive_comb_survey} for more details.\n\nThe most relevant Borel graphs for us are the ones that are induced by a free Borel action of $\\mathbb{Z}^d$.\nThat is, there are $d$-many aperiodic Borel bijections $\\varphi_i:X\\to X$ that pairwise commute, where \\emph{aperiodic} means that $\\varphi_i^k(x)\\not=x$ for every $k>0$, $i\\in [d]$ and $x\\in X$.\nThen the Borel (labeled) graph $\\mathcal{G}$ that is induced by this action consists of edges of the form $(x,\\varphi_i(x))$, where $x\\in X$ and $i\\in [d]$.\nThe most important result in this context is that $\\mathcal{G}$ admits a \\emph{Borel toast structure}.\nThis follows from work of Gao et al \\cite{GJKS2}, see Marks and Unger \\cite{Circle} for proof.\nWe note that it is not understood in general what Borel graphs admit a Borel toast structure.\nThis is connected with the notoriously difficult question about Borel hyperfiniteness, see \\cite{AsymptoticDim}.\n\nGiven an LCL $\\Pi$, we say that $\\Pi$ is in the class $\\mathsf{BOREL}$ if there is a Borel measurable $\\Pi$-coloring of every Borel graph $\\mathcal{G}$ as above, i.e., induced by a free Borel action of $\\mathbb{Z}^d$.\nEquivalently, it is enough to decide whether $\\Pi$ can be solved on the canonical Borel graph $\\mathcal{G}$ on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$.\nMoreover, if $\\Pi$ has some inputs $\\Sigma_{in}$, then we consider any Borel labeling $X\\to \\Sigma_{in}$ and then ask if there is a Borel solution that satisfies all the constraints on $\\Sigma_{in}\\times \\Sigma_{out}$.\n\n\\begin{theorem}[Basically \\cite{GJKS}]\nLet $\\Pi$ be an LCL that is in the class $\\LTOAST$.\nThen $\\Pi\\in\\mathsf{BOREL}$.\n\\end{theorem}\n\nIn fact, all the LCL problems without input that we are aware of and that are in the class $\\mathsf{BOREL}$ are in fact in the class $\\TOAST$.\nIt is an open problem to decide whether it is always the case.\n\nAnother interesting perspective how to see Borel constructions is as a countably infinite analogue of $\\mathsf{LOCAL}(O(\\log^*(n)))$.\nNamely, suppose that you are allowed to inductively (a) construct Borel MIS with some parameter (b) do local $O(1)$-construction that depends on the previously build structure.\nIt is clear that objects that are build in this way are Borel, however, it is not clear if this uses the whole potential of Borel sets to find Borel solutions of LCLs.\n\nAs a last remark we note that traditionally in descriptive combinatorics, people consider relaxations of Borel constructions.\nThis is usually done via ignoring how a solution behave on some negligible Borel set.\nProminent examples are \\emph{measure $0$ sets} for some Borel probability measure on $(X,\\mathcal{B})$ or \\emph{meager sets} for some compatible complete metric on $(X,\\mathcal{B})$.\nWe denote the classes of problems that admit such a solution for \\emph{every} Borel probability measure or compatible complete metric as $\\mathsf{MEASURE}$ and $\\mathsf{BAIRE}$, respectively.\nWithout going into any details we mention that on general graphs these classes are much richer than $\\mathsf{BOREL}$, see e.g. \\cite{DetMarks,BrooksMeas}.\nHowever, it is still plausible that on $\\mathbb{Z}^d$ we have $\\mathsf{BOREL}=\\mathsf{BAIRE}=\\mathsf{MEASURE}$.\nIt is well-know that this is true in $\\mathbb{Z}$ \\cite{grebik_rozhon2021LCL_on_paths}.\n\n\n\\begin{comment}\nIn both cases we are allowed to neglect some small sets, where we either do not care about the solution or are allowed to use the axiom of choice.\nSuppose that $\\mathcal{G}$ is a Borel graph on a standard Borel space $(X,\\mathcal{B})$ and $\\Pi$ is an LCL.\nSuppose, in addition, that there is a Borel probability measure $\\mu$ on $(X,\\mathcal{B})$, that is a $\\sigma$-additive measure on $\\mathcal{B}$, see \\cite[Section~?]{Kec}.\nWe say that $\\Pi$ admits \n\n\n\n\n\n\n\n\n\n\nand show the arrows that are left to complete Figure~\\ref{fig:big_picture_grids}.\nRecall that a \\emph{Borel graph} $\\mathcal{G}$ is a triplet $\\mathcal{G}=(X,\\mathcal{B},E)$, where $(X,\\mathcal{B})$ is a standard Borel space, $(X,E)$ is a graph and $E\\subseteq X\\times X$ is an element of the product $\\sigma$-algebra.\nFor the purposes of this paper we always assume that $\\mathcal{G}$ is induced by a free Borel action of $\\mathbb{Z}^d$ for some $d>0$, i.e., there is a Borel action $\\mathbb{Z}^d\\curvearrowright X$ with trivial stabilizers such that $\\mathcal{G}$ is induced by the standard generators of $\\mathbb{Z}^d$.\n\nLet $\\Pi$ be an LCL.\nThen $\\Pi$ belongs to the class $\\mathsf{BOREL}$ if there is a Borel $\\Pi$-coloring of every such Borel graph $\\mathcal{G}$.\nMoreover, we say that $\\Pi$ is in the class $\\mathsf{MEASURE}$ if for every such Borel graph $\\mathcal{G}$ and for every \\emph{Borel probability measure} $\\mu$ on $(X,\\mathcal{B})$ there is a Borel $\\Pi$-coloring of $\\mathcal{G}$ on a $\\mu$-conull set.\nSimilarly, we say that $\\Pi$ is in the class $\\mathsf{BAIRE}$ if for every such Borel graph $\\mathcal{G}$ and for every \\emph{compatible Polish topology} $\\tau$ on $(X,\\mathcal{B})$ there is a Borel $\\Pi$-coloring of $\\mathcal{G}$ on a $\\tau$-comeager set.\nIt follows directly from the definition that $\\mathsf{BOREL}\\subseteq \\mathsf{MEASURE},\\mathsf{BAIRE}$.\n\\end{comment}\n\n\n\n\\begin{comment}\nNext we follow Gao et al. \\cite{GJKS}.\nDenote as $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ the free part of the Bernoulli shift $\\mathbb{Z}^d\\curvearrowright 2^{\\mathbb{Z}^d}$ endowed with the inherited product topology.\nThen it is easy to see that $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ is a Polish space and that the action is continuous.\nConsider the Borel graph $\\mathcal{G}$ on $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ induced by the standard set of generators.\nThis graph is universal in the sense that an LCL $\\Pi$ is in the class $\\mathsf{BOREL}$ if and only if there is a Borel $\\Pi$-coloring of $\\mathcal{G}$, this follows from result of Seward and Tucker-Drob~\\cite{?}.\nWe say that an LCL $\\Pi$ is in the class $\\mathsf{CONTINUOUS}$ if there is a $\\Pi$-coloring of $\\mathcal{G}$ that is continuous.\nThis is the same as saying that that there is a continuous equivariant map from $\\operatorname{Free}(2^{\\mathbb{Z}^d})$ to $X_\\Pi$.\nThe following is a corollary of the Twelve Tiles Theorem of Gao et al. \\cite{GJKS}.\n\n\\begin{theorem}\\cite[Theorem 5.5]{GJKS}\n$\\mathsf{CONTINUOUS}=\\mathsf{LOCAL}(O(\\log^* n))$.\n\\end{theorem}\n\nWe note that in the general case of all groups this equality was shown by Bernshteyn\\cite{?} and, independently, by Seward. \n\nIt remains to describe how toastable algorithms give Borel solutions.\nAs an immediate corollary we get that all the problems from Section~\\ref{sec:SpecificProblems} are in the class $\\mathsf{BOREL}$.\n\n\\begin{theorem}\n$\\LTOAST \\subseteq \\mathsf{BOREL}$.\n\\end{theorem}\n\\begin{proof}\nSuppose that $\\Pi$ is in the class $\\LTOAST$ and let $f$ be an toast algorithm with parameter $k$.\nIt follows from work of Gao et al \\cite{GJKS2}, see Marks and Unger \\cite{Circle} for proof, that every Borel graph $\\mathcal{G}$ (that is induced by a free action of $\\mathbb{Z}^d$) admits a Borel $k$-toast for every $k\\in \\mathbb{N}$.\nIt is a standard fact that using $f$ inductively on the depth of the toast produces a Borel $\\Pi$-coloring of $\\mathcal{G}$.\n\\end{proof}\n\\end{comment}\n\n\n\n\n\n\\subsection{Finitely Dependent Processes}\n\nIn this section we use the notation that was developed in \\cref{subsec:subshifts}, also we restrict our attention to LCLs without input.\nLet $\\Pi=(b,t,\\mathcal{P})$ be such an LCL.\nRecall that $X_\\Pi$ is the space of all $\\Pi$-colorings.\n\nWe say that a distribution $\\mu$ on $b^{\\mathbb{Z}^d}$ is \\emph{finitely dependent} if it is shift-invariant and there is $k\\in \\mathbb{N}$ such that the events $\\mu\\upharpoonright A$ and $\\mu \\upharpoonright B$ are independent for every finite set $A,B\\subseteq \\mathbb{Z}^d$ such that $d(A,B)\\ge k$.\nNote that shift-invariance is the same as saying that the distribution $\\mu\\upharpoonright A$ does not depend on the position of $A$ in $\\mathbb{Z}^d$ but rather on its shape.\n\nHere we study the following situation:\nSuppose that there is a finitely dependent process $\\mu$ that is concentrated on $X_\\Pi$.\nIn that case we say that $\\Pi$ is in the class $\\mathsf{FINDEP}$. \nIn other words, $\\Pi \\in \\mathsf{FINDEP}$ if there is a distribution from which we can sample solutions of $\\Pi$ such that if we restrict the distribution to any two sufficient far sets, the corresponding restrictions of the solution of $\\Pi$ that we sample are independent. \n\nIt is easy to see that all problems in $\\mathsf{LOCAL}(O(1))$ (that is, problems solvable by a block factor) are finitely dependent, but examples of finitely dependent processes outside of this class were rare and unnatural, see the discussion in \\cite{FinDepCol}, until Holroyd and Liggett \\cite{FinDepCol} provided an elegant construction for 3- and 4-coloring of a path and, more generally, for distance coloring of a grid with constantly many colors. This implies the following result.\n\n\\begin{theorem}[\\cite{FinDepCol}]\n$\\mathsf{LOCAL}(O(\\log^* n)) \\subseteq \\mathsf{FINDEP}$ (see \\cref{fig:big_picture_grids}).\n\\end{theorem} \n\nOn the other hand, recently Spinka \\cite{Spinka} proved that any finitely dependent process can be described as a finitary factor of iid process.\\footnote{His result works in a bigger generality for finitely generated amenable groups.}\nUsing \\cref{thm:toast_in_tail}, we sketch how this can be done, we keep the language of LCLs.\n\n\\begin{theorem}[\\cite{Spinka}]\nLet $\\Pi$ be an LCL on $\\mathbb{Z}^d$ and $\\mu$ be finitely dependent process on $X_\\Pi$, i.e., $\\Pi\\in \\mathsf{FINDEP}$.\nThen there is a finitary factor of iid $F$ (equivalently an oblivious local algorithm $\\mathcal{A}$) with tail decay $1\/r^{1- O(\\sqrt{\\log r})}$ such that the distribution induced by $F$ is equal to $\\mu$.\nIn fact, $\\Pi$ is in the class $\\TOAST$. \n\\end{theorem}\n\\begin{proof}[Proof Sketch]\nSuppose that every vertex is assigned uniformly two reals from $[0,1]$.\nUse the first string to produce a $k$-toast as in \\cref{thm:toast moment}, where $k$ is from the definition of $\\mu$.\n\nHaving this use the second random string to sample from $\\mu$ inductively along the constructed $k$-toast.\nIt is easy to see that this produces ffiid $F$ that induces back the distribution $\\mu$.\nIn fact, it could be stated separately that the second part describes a randomized $\\TOAST$ algorithm: always samples from $\\mu$ given the coloring on smaller pieces that was already fixed. The finitely dependent condition is exactly saying that the joint distribution over the solution on the small pieces of $\\mu$ is the same as the distribution we sampled from.\nNow use \\cref{thm:derandomOfToast} to show that $\\Pi\\in \\TOAST$. \n\\end{proof}\n\nIn our formalism, $\\mathsf{FINDEP} \\subseteq \\TOAST$, or equivalently we get $\\mathsf{TAIL}(\\mathrm{tower}(\\Omega(r))) \\subseteq \\mathsf{FINDEP} \\subseteq \\mathsf{TAIL}(r^{1-o(1)})$.\nIt might be tempting to thing that perhaps $\\mathsf{FINDEP} =\\TOAST$.\nIt is however a result of Holroyd et al \\cite[Corollary~25]{HolroydSchrammWilson2017FinitaryColoring} that proper vertex $3$-coloring, $\\Pi^{d}_{3,\\chi}$, is not in $\\mathsf{FINDEP}$ for every $d>1$.\nThe exact relationship between these classes remain mysterious.\n\n\n\n\\begin{comment}\n\\begin{theorem}\nLet $\\Pi$ be an LCL on $\\mathbb{Z}^d$.\nSuppose that $\\Pi$ is in the class $\\mathsf{TAIL}(\\tow(\\Omega(r)))$.\nThen there is finitely dependent process $\\mu$ on $X_\\Pi$.\n\\end{theorem}\n\\begin{proof}\nIt follows from Theorem~\\ref{thm:classification_of_local_problems_on_grids} that there is $r,\\ell \\in \\mathbb{N}$ such that $\\ell\\le r-t$ and a rule $\\mathcal{A}$ of locality at most $\\ell$ that produces a $\\Pi$-coloring for every $(r+t)$-power proper vertex coloring of $\\mathbb{Z}^d$.\n\nSuppose that $\\mu_{(r+t)}$ is a finitely dependent $(r+t)$-power proper vertex coloring of $\\mathbb{Z}^d$ process with locality $k\\in \\mathbb{N}$.\nWe claim that composing $\\mu_{(r+t)}$ with $\\mathcal{A}$ yields a finitely dependent process with locality $(2\\ell+k)$.\nLet $A\\subseteq \\mathbb{Z}^d$ be a finite set.\nThen the output of $\\mathcal{A}$ depends on $\\mathcal{B}(A,\\ell)$.\nTherefore if $A,B\\subseteq \\mathbb{Z}^d$ satisfy $d(A,B)>2\\ell+k$, then\n$$d(\\mathcal{B}(A,\\ell),\\mathcal{B}(B,\\ell))>k$$\nand we conclude that the outputs of $\\mathcal{A}$ on $A$ and $B$ are independent.\n\nIt remains to find $\\mu_{r}$ finitely dependent $r$-power proper vertex coloring of $\\mathbb{Z}^d$ process for every $r>0$.\nThis is done by Holroyd and Liggett \\cite[Theorem~20]{FinDepCol}.\n{\\color{red} or something like that, they use a lot of colors fr $r$-power coloring...}\n\\end{proof}\n\nGiven a fiid $F$ we define $F^*\\lambda$ to be the push-forward of the product measure via $F$.\nSince $F$ is equivariant it follows that $F^*\\lambda$ is shift-invariant.\n\nThe following theorem was essentially proven by Spinka\\cite{?}. His result nicely fits to our definitions in \\cref{sec:Toast}. \n\\begin{theorem}\n\n\\end{theorem}\n\n\n\\begin{theorem}\nLet $\\Pi$ be an LCL problem on $\\mathbb{Z}^d$ and $\\mu$ be finitely dependent process on $X_\\Pi$.\nThen for every $\\epsilon>0$ there is a ffiid $F$ such that $F^*\\lambda=\\mu$ and the tail of $F$ has the $(1-\\epsilon)$-moment finite.\nMoreover, $\\Pi$ is in the class $\\TOAST$.\n\\end{theorem}\n\n\\begin{remark}\nMaybe the random toast construction works for any finitely dependent processes not just LCL.\n\nThere is the notion of strongly irreducible shifts of finite type, some comment about it...\n\nDefine $\\mathsf{FINDEP}$ to be the class of all LCL problems that admit a finitely dependent distribution $\\mu$.\nIt follows that $\\mathsf{TAIL}(\\tow(\\Omega(r)))\\subseteq \\mathsf{FINDEP} \\subseteq \\mathsf{TOAST}$ and we have $\\mathsf{FINDEP}\\not = \\TOAST$ because proper vertex $3$-coloring $\\Pi^d_{\\chi,3}$ is not in $\\mathsf{FINDEP}$ for every $d>1$, by \\cite[Corollary~25]{HolroydSchrammWilson2017FinitaryColoring}.\n\\end{remark}\n\n\n\\begin{proof}\nWe define a spontaneous randomized toast algorithm $\\Omega$ such that $\\Omega(\\mathcal{D},\\ell)$ is a $\\Pi$-coloring for every $k$-toast $\\mathcal{D}$ and every $\\ell\\in [0,1]^{\\mathbb{Z}^d}$.\nThis already implies the moreover part by Theorem~\\ref{thm:derandomOfToast}.\n\nLet $A\\subseteq \\mathbb{Z}^d$ be a finite subset of $\\mathbb{Z}^d$.\nWrite $\\mathbb{P}_A$ for the distribution of $\\Pi$-patterns of $\\mu$ on $A$, i.e., $\\mathbb{P}_A=\\mu\\upharpoonright A$.\nSuppose now that $(A,c,\\ell)$ are given, where $A\\subseteq \\mathbb{Z}^d$ is finite subset, $c$ is a partial coloring of $A$ and $\\ell\\in [0,1]^{A}$.\nConsider the conditional distribution $\\mathbb{P}_{A,c}$ of $\\mathbb{P}_A$ given the event $c$.\nIf $\\mathbb{P}_{A,c}$ is non-trivial, then define $\\Omega$ such that $\\Omega(A,c,{-})$ induces the same distribution as $\\mathbb{P}_{A,c}$ when $\\ell$ is fixed on $\\operatorname{dom}(c)$.\nOtherwise define $\\Omega$ arbitrarily.\n\nLet $\\mathcal{D}$ be a $k$-toast.\nWe show that $\\Omega(\\mathcal{D},\\ell)$ is a $\\Pi$-coloring for every $\\ell\\in [0,1]^{\\mathbb{Z}^d}$ and $\\Omega(\\mathcal{D},{-})$ produces $\\mu$.\nNote that the first claim shows that $\\Pi$ is in the class $\\RTOAST$, while the second produces $\\mu$ as ffiid as follows.\nSuppose that every $v\\in \\mathbb{Z}^d$ is assigned two independent reals $l_{v,0}$ and $l_{v,1}$.\nUse the first real to produce a $k$-toast as in Theorem~\\ref{thm:toast moment} and use the second real together with the produced toast as an input to $\\Omega$.\nIt follows from the independence of $l_{v,0}$ and $l_{v,1}$ that this produces $\\mu$.\n\nWe show that\n$$\\Omega(\\mathcal{D},{-})\\upharpoonright D=\\mu\\upharpoonright D$$\nfor every $D\\in \\mathcal{D}$ inductively on $\\mathfrak{r}_{\\mathcal{D}}$.\nIf $\\mathfrak{r}_{\\mathcal{D}}(D)=0$, then it follows from the definition.\nSuppose that $\\mathfrak{r}_{\\mathcal{D}}(D)>0$ and let $E_1,\\dots, E_t\\in \\mathcal{D}$ be an enumeration of $E\\subset D$.\nIt follows from the inductive hypothesis together with the assumption that $\\mu$ is $k$-independent and $\\mathcal{D}$ is a $k$-toast that \n$$\\Omega(\\mathcal{D},{-})\\upharpoonright \\bigcup_{i\\le t}E_i=\\mu\\upharpoonright \\bigcup_{i\\le t} E_i.$$\nNow when $\\ell$, and therefore a partial coloring $c$, is fixed on $\\bigcup_{i\\le t}E_i=\\operatorname{dom}(c)$, the definition $\\Omega(D,c,{-})$ extends $c$ in such a way that $\\mathbb{P}_{D,c}=\\mu\\upharpoonright \\operatorname{dom}(c)$.\nSince this holds for any such assignment $\\ell$ of real numbers on $\\bigcup_{i\\le t}E_i$ we deduce that\n$$\\Omega(\\mathcal{D},{-})\\upharpoonright D=\\mu\\upharpoonright D$$\nand the proof is finished.\n\\end{proof}\n\\end{comment}\n\n\n\n\n\n\n\\subsection{$\\mathsf{TAIL}$ Complexities in Distributed and Centralised Models}\n\\label{subsec:congestLCA}\n\nAlthough we are not aware of a systematic work on $\\mathsf{TAIL}$ complexities in the distributed algorithms community, similar notions like oblivious algorithms \\cite{Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring}, node-average complexity \\cite{feuilloley2017node_avg_complexity_of_col_is_const,barenboim2018avg_complexity,chatterjee2020sleeping_MIS} or energy complexity \\cite{chang2018energy_complexity_exp_separation,chang2020energy_complexity_bfs} are studied. \nThe breakthrough MIS algorithm of Ghaffari \\cite{ghaffari2016MIS} may have found many applications exactly because of the fact that its basis is an algorithm with exponential tail behavior. \n\nWe discuss two applications of algorithms with nontrivial $\\mathsf{TAIL}$ complexities. First, the node-average complexity literature studies the average coding radius that a node needs to answer. Whenever an oblivious algorithm has its first moment finite, we have an algorithm with constant average coding radius. \nWith a little more work, one can see that the result from \\cref{thm:3in toast}, that is $3$-coloring of grids has $\\mathsf{TAIL}$ complexity $r^{1-O(1\/\\sqrt{\\log r})}$ can be adapted to the setting of finite grids where it gives a $3$-coloring algorithm with average complexity $2^{O(\\sqrt{\\log n})}$. \nFor nonconstant $\\Delta$, the situation is more complicated. In \\cite{barenboim2018avg_complexity} it was shown that the Luby's $\\Delta+1$ coloring algorithm has tail $2^{\\Omega(r)}$ even for nonconstant $\\Delta$, but the question of whether MIS can have constant average coding radius in nonconstant degree graphs is open \\cite{chatterjee2020sleeping_MIS}. \n\nAnother application of the $\\mathsf{TAIL}$ bound approach comes from the popular centralised $\\mathsf{LCA}$ model \\cite{lca11,lca12}: there one has a huge graph and instead of solving some problem, e.g., coloring, on it directly, we are only required to answer queries about the color of some given nodes. It is known that, as in the $\\mathsf{LOCAL}$ model, one can solve $\\Delta+1$ coloring with $O(\\log^* n)$ queries on constant degree graphs by a variant of Linial's algorithm \\cite{even03_conflictfree}. But using the tail version of the Linial's algorithm instead, we automatically get that $O(\\log^* n)$ is only a worst-case complexity guarantee and on average we expect only $O(1)$ queries until we know the color of a given node. That is, the amortized complexity of a query is in fact constant. The same holds e.g. for the $O(\\log n)$ query algorithm for LLL from \\cite{brandt_grunau_rozhon2021LLL_in_LCA}. \n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Open Questions}\n\\label{sec:open_problems}\n\nHere we collect some open problems that we find interesting. \n\n\\begin{enumerate}\n \\item Is it true that $\\RLTOAST = \\LTOAST$? If yes, this implies $\\mathsf{TAIL} \\subseteq \\mathsf{BOREL}$. \n \\item Is it true that $\\RLTOAST = \\RTOAST$? If yes, this implies $\\mathsf{TAIL} = \\mathsf{TAIL}(r^{1-o(1)})$ and answers the fourth question of Holroyd et al. \\cite{HolroydSchrammWilson2017FinitaryColoring}. \n\t\\item Can the tail complexity $r \/ 2^{O(\\sqrt{\\log r})}$ in the construction of toast be improved to $r \/ \\operatorname{poly} \\log (r)$? The same can be asked for the complexity of intermediate problems from \\cref{sec:Inter}. \n\t\t\\item Is there a local coloring problem $\\Pi$ such that ${\\bf M}(\\Pi)$ is non-integer value or $0$, i.e., ${\\bf M}(\\Pi)\\in [0,d-1]\\setminus \\{1,2,\\dots, d-1\\}$?\n\t\t\\item Can you find an example of a local problem without input labels which does not belong to any class in \\cref{fig:big_picture_grids} but the solutions to this problem are not ``periodic'' (this cannot happen for $d=1$). Maybe (a version of) Penrose tilings can be helpful there?\n\t\\item What is the value ${\\bf M}(\\Pi^d_{pm})$, where $\\Pi^d_{pm}$ is the problem of finding a perfect matching of $\\mathbb{Z}^d$, where $d>2$?\n \\item Let $d>1$.\n\tFor what finite collections of $d$-dimensional rectangles $\\mathcal{R}$ is $\\Pi^d_{\\mathcal{R}}$ is in the class $\\mathsf{LOCAL}(O(\\log^* n))$? Is there a finite collection of $d$-dimensional rectangles $\\mathcal{R}$ that would satisfy $1< {\\bf M}(\\Pi^d_{\\mathcal{R}})\\le d-1$?\n\t\\item Is it true (on $\\mathbb{Z}^d$) that $\\mathsf{BOREL}=\\mathsf{MEASURE}$ or $\\mathsf{BOREL} = \\mathsf{BAIRE}$?\n\\end{enumerate}\n\n\\section*{Acknowledgement}\nWe would like to thank Ferenc Bencs, Anton Bernshteyn, Sebastian Brandt, Yi-Jun Chang, Mohsen Ghaffari, Aranka Hru\u0161kov\u00e1, Stephen Jackson, Oleg Pikhurko, Brandon Sewart, Yinon Spinka, L\u00e1szlo T\u00f3th, and Zolt\u00e1n Vidny\u00e1szky for many engaging discussions. \n\nThe first author was supported by Leverhulme Research Project Grant RPG-2018-424. \nThis project has received funding from the European Research Council (ERC) under the European\nUnions Horizon 2020 research and innovation programme (grant agreement No. 853109). \n\n\\bibliographystyle{alpha}\n\n\n\\section{Introduction}\n\\label{sec:intro}\nIn this work, we study local problems on regular trees from three different perspectives. \n\n\\textit{First,} we consider the perspective of distributed algorithms. In distributed computing, the studied setup is a network of computers where each computer can only communicate with its neighbors. Roughly speaking, the question of interest in this area is which problems can be solved with only a few rounds of communication in the underlying network. \n\n\\textit{Second}, we consider the perspective of (finitary) factors of iid processes.\nIn probability, random processes model systems that appear to vary in a random manner.\nThese include Bernoulli processes, Random walks etc.\nA particular, well-studied, example is the Ising model.\n\n\\textit{Third}, we investigate the perspective of descriptive combinatorics. The goal of this area is to understand which constructions on infinite graphs can be performed without using the so-called axiom of choice.\n\nAlthough many of the questions of interest asked in these three areas are quite similar to each other, no systematic connections were known until an insightful paper of Bernshteyn~\\cite{Bernshteyn2021LLL} who showed that results from distributed computing can automatically imply results in descriptive combinatorics.\nIn this work, we show that the connections between the three areas run much deeper than previously known, both in terms of techniques and in terms of complexity classes.\nIn fact, our work suggests that it is quite useful to consider all three perspectives as part of a common theory, and we will attempt to present our results accordingly.\nWe refer the reader to \\cref{fig:big_picture_trees} for a partial overview of the rich connections between the three perspectives, some of which are proven in this paper. \n\nIn this work, we focus on the case where the graph under consideration is a regular tree.\nDespite its simplistic appearance, regular trees play an important role in each of the three areas, as we will substantiate at the end of this section.\nTo already provide an example, in the area of distributed algorithms, almost all locality lower bounds are achieved on regular trees.\nMoreover, when regarding lower bounds, the property that they already apply on regular trees actually \\emph{strengthens} the result---a fact that is quite relevant for our work as our main contribution regarding the transfer of techniques between the areas is a new lower bound technique in the area of distributed computation that is an adaptation and generalization of a technique from descriptive combinatorics.\nRegarding our results about the relations between complexity classes from the three areas, we note that such connections are also studied in the context of paths and grids in two parallel papers \\cite{grebik_rozhon2021LCL_on_paths,grebik_rozhon2021toasts_and_tails}.\n\nIn the remainder of this section, we give a high-level overview of the three areas that we study.\nThe purpose of these overviews is to provide the reader with a comprehensive picture of the studied settings that can also serve as a starting point for delving deeper into selected topics in those areas---in order to follow our paper, it is not necessary to obtain a detailed understanding of the results and connections presented in the overviews. Necessary technical details will be provided in \\cref{sec:preliminaries}.\nMoreover, in \\cref{sec:contribution}, we present our contributions in detail.\n\n\n\n\n\n\\paragraph{Distributed Computing}\nThe definition of the $\\mathsf{LOCAL}$ model of distributed computing by Linial~\\cite{linial92LOCAL} was motivated by the desire to understand distributed algorithms in huge networks. \nAs an example, consider a huge network of wifi routers. Let us think of two routers as connected by an edge if they are close enough to exchange messages. \nIt is desirable that such close-by routers communicate with user devices on different channels to avoid interference.\nIn graph-theoretic language, we want to properly color the underlying network. \nEven if we are allowed a color palette with $\\Delta+1$ colors where $\\Delta$ denotes the maximum degree of the graph (which would admit a simple greedy algorithm in a sequential setting), the problem remains highly interesting in the distributed setting, as, ideally, each vertex decides on its output color after only a few rounds of communication with its neighbors, which does not allow for a simple greedy solution. \n\nThe $\\mathsf{LOCAL}$ model of distributed computing formalizes this setup: we have a large network, where each vertex knows the network's size, $n$, and perhaps some other parameters like the maximum degree $\\Delta$. In the case of randomized algorithms, each vertex has access to a private random bit string, while in the case of deterministic algorithms, each vertex is equipped with a unique identifier from a range polynomial in the size $n$ of the network. \nIn one round, each vertex can exchange any message with its neighbors and can perform an arbitrary computation. The goal is to find a solution to a given problem in as few communication rounds as possible. \nAs the allowed message size is unbounded, a $t$-round $\\mathsf{LOCAL}$ algorithm can be equivalently described as a function that maps $t$-hop neighborhoods to outputs---the output of a vertex is then simply the output its $t$-hop neighborhood\nmapped to by this function.\nAn algorithm is correct if and only if the collection of outputs at all vertices constitutes a correct solution to the problem. \n\nThere is a rich theory of distributed algorithms and the local complexity of many problems is understood. \nThe case of trees is a highlight of the theory: it is known that any local problem (a class of natural problems we will define later) belongs to one of only very few complexity classes.\nMore precisely, for any local problem, its randomized local complexity is either $O(1), \\Theta(\\log^* n), \\Theta(\\log\\log n), \\Theta(\\log n)$, or $\\Theta(n^{1\/k})$ for some $k \\in \\mathbb{N}$. Moreover, the deterministic complexity is always the same as the randomized one, except for the case $\\Theta(\\log\\log n)$, for which the corresponding deterministic complexity is $\\Theta(\\log n)$ (see \\cref{fig:big_picture_trees}). \n\n\n\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=\\textwidth]{fig\/big_picture_trees.pdf}\n \\caption{Complexity classes on regular trees considered in the three areas of distributed computing, factors of iid processes\/uniform algorithms, and descriptive combinatorics. The left part shows complexity classes of distributed computing. We use the shorthand $\\mathsf{LOCAL}$ if it does not matter whether we talk about the deterministic or randomized complexity. These two notions differ only for the class of problems of randomized local complexity $O(\\log\\log n)$, which have deterministic complexity $O(\\log n)$. \\\\\n The uniform complexity classes of sublogarithmic complexity are in correspondence to appropriate classes in the randomized local complexity model, as proven in \\cref{sec:separating_examples}. On the other hand, the class fiid is very similar to the class $\\mathsf{MEASURE}$ from descriptive combinatorics. The equivalence of the class $\\mathsf{CONTINUOUS}$ and $\\mathsf{LOCAL}(O(\\log^* n)) = \\mathsf{OLOCAL}(O(\\log^* 1\/\\varepsilon))$ is marked with a dashed arrow as it was proven in case the tree is generated by a group action (think of the tree being equipped with an additional $\\Delta$-edge coloring). The inclusion $\\mathsf{LOCAL}(O(\\log^* n)) \\subseteq \\mathsf{BOREL}$ however clearly holds also in our setting. The class $\\mathsf{BOREL}$ is incomparable with $\\mathsf{RLOCAL}(O(\\log\\log n))$, as proven in \\cref{sec:separating_examples}. }\n \\label{fig:big_picture_trees}\n\\end{figure}\n\n\n\\todo{thm 5.1 - 5.3}\n\n\n\n\n\\paragraph{(Finitary) Factors of iid Processes and Uniform Algorithms}\n\nIn recent years, \\emph{factors of iid (fiid) processes} on trees attracted a lot of attention in combinatorics, probability, ergodic theory and statistical physics \\cite{ABGGMP,ArankaFeriLaci,BackSzeg,BackVir,backhausz,backhausz2,BGHV,Ball,Bowen1,Bowen2,CPRR,csoka2015invariant,GabLyo,GamarnikSudan,GGP,HarangiVirag,HolroydPoisson,HLS,HPPS,HolroydPeres,HoppenWormald,Kun,RahmanPercolation,RahVir,LyonsTrees,Mester,schmidt2016circular,Timar1,Timar2}.\nIntuitively, factors of iid processes are randomized algorithms on, e.g., infinite $\\Delta$-regular trees, where each vertex outputs a solution to a problem after it explores random strings on vertices of the whole tree.\nAs an example, consider the \\emph{perfect matching} problem.\nAn easy parity argument shows that perfect matching cannot be solved by any local randomized algorithm on finite trees.\nHowever, if we allow a small fraction of vertices not to be matched, then, by a result of Nguyen and Onak \\cite{nguyen_onak2008matching} (see also \\cite{ElekLippner}), there is a constant-round randomized algorithm that produces such a matching on high-girth graphs (where the constant depends on the fraction of unmatched vertices that we allow).\nThis result can also be deduced from a result of Lyons and Nazarov \\cite{lyons2011perfect}, who showed that perfect matching can be described as a factor of iid process on an infinite $\\Delta$-regular tree.\nThe high-level idea behind this connection is that high-girth graphs approximate the infinite $\\Delta$-regular tree and constant-round local algorithms approximate factors of iid processes.\nThis correspondence is formalized in the notion of \\emph{Benjamini-Schramm} or \\emph{local-global} convergence \\cite{benjamini2011recurrence, hatamilovaszszegedy}.\nWe note that getting only ``approximate'' solutions, that is, solutions where a small fraction of vertices does not have to satisfy the constraints of a given problem, is intrinsic in this correspondence.\nRegardless, there are many techniques, such as entropy inequality \\cite{backhausz} or correlation decay \\cite{backhausz2}, and particular results such as the aforementioned perfect matching problem \\cite{lyons2011perfect} that provide lower and upper bounds, respectively, in our setting as well.\nWe refer the reader to \\cite{LyonsTrees,BGHV} for a comprehensive summary of the field. \n\n\nIn this paper, we mostly consider a stronger condition than fiid, namely so-called \\emph{finitary} factors of iid (ffiid) processes that are studied in the same context as fiid \\cite{holroyd2017one_dependent_coloring,HolroydSchrammWilson2017FinitaryColoring,Spinka}. \nPerhaps surprisingly, the notion of ffiid is identical to the notion of so-called uniform distributed randomized algorithms \\cite{Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring,grebik_rozhon2021toasts_and_tails} that we now describe. \nWe define an uniform local algorithm as a randomized local algorithm that does not know the size of the graph $n$ -- this enables us to run such an algorithm on infinite graphs, where there is no $n$. \nMore precisely, we require that each vertex eventually outputs a solution that is compatible with the output in its neighborhood, but the time until the vertex finishes is a potentially unbounded random variable. \nAs in the case of classical randomized algorithms, we can now measure the \\emph{uniform complexity} of an uniform local algorithm (known as the tail decay of ffiid \\cite{HolroydSchrammWilson2017FinitaryColoring}). \nThe uniform complexity of an algorithm is defined as the function $t(\\varepsilon)$ such that the probability that the algorithm run on a specific vertex needs to see outside its $t(\\varepsilon)$-hop neighborhood is at most $\\varepsilon$. \nAs in the case of classical local complexity, there is a whole hierarchy of possible uniform complexities (see \\cref{fig:big_picture_trees}). \n\nWe remark that uniform distributed local algorithms can be regarded as Las Vegas algorithms. The output will always be correct; there is however no fixed guarantee at what point all vertices have computed their final output. On the other hand, a randomized distributed local algorithm can be viewed as a Monte Carlo algorithm as it needs to produce an output after a fixed number of rounds, though the produced output might be incorrect. \n\n\n\n\n\n\n\n\n\\paragraph{Descriptive Combinatorics}\n \n \n\t\n\tThe Banach-Tarski paradox states that a three-dimensional ball of unit volume can be decomposed into finitely many pieces that can be moved by isometries (distance preserving transformations such as rotations and translations) to form two three-dimensional balls each of them with unit volume(!).\n The graph theoretic problem lurking behind this paradox is the following: fix finitely many isometries of $\\mathbb{R}^3$ and then consider a graph where $x$ and $y$ are connected if there is an isometry that sends $x$ to $y$.\n\tThen our task becomes to find a perfect matching in the appropriate subgraph of this graph -- namely, the bipartite subgraph where one partition contains points of the first ball and the other contains points of the other two balls. \n\tBanach and Tarski have shown that, with a suitably chosen set of isometries, the axiom of choice implies the existence of such a matching. \n\tIn contrast, since isometries preserve the Lebesgue measure, the pieces in the decomposition cannot be Lebesgue measurable. \n\tSurprisingly, Dougherty and Foreman \\cite{doughertyforeman} proved that the pieces in the Banach-Tarski paradox can have the \\emph{Baire property}. \n\tThe Baire property is a topological analogue of being Lebesgue measurable; a subset of $\\mathbb{R}^3$ is said to have the Baire property if its difference from some open set is topologically negligible.\n\t\n \n\t \n\t\n\t\n Recently, results similar to the Banach-Tarski paradox that lie on the border of combinatorics, logic, group theory, and ergodic theory led to an emergence of a new field often called \\textit{descriptive} or \\textit{measurable combinatorics}. \n The field focuses on the connection between the discrete and continuous and is largely concerned with the investigation of graph-theoretic concepts. \n The usual setup in descriptive combinatorics is that we have a graph with uncountably many connected components, each being a countable graph of bounded degree. \n For example, in case of the Banach-Tarski paradox, the vertices of the underlying graph are the points of the three balls, edges correspond to isometries, and the degree of each vertex is bounded by the number of chosen isometries. \n Some of the most beautiful results related to the field include \\cite{laczk, marksunger,measurablesquare,doughertyforeman,marks2016baire,gaboriau,KST,DetMarks,millerreducibility,conley2020borel,csokagrabowski,Bernshteyn2021LLL}, see \\cite{kechris_marks2016descriptive_comb_survey,pikhurko2021descriptive_comb_survey} for recent surveys. \n \n Importantly, in many of these results, including the Banach-Tarski paradox, graphs where each component is an infinite $\\Delta$-regular tree appear naturally. \n Oftentimes, questions considered in descriptive combinatorics lead to constructing a solution to a local problem in the underlying uncountable graph (in the case of Banach-Tarski, the local problem is perfect matching). \n The construction needs to be such that the solution of the problem has some additional regularity properties. For example in the case of Banach-Tarski, a solution is possible when the regularity condition is the Baire property, but not if it is Lebesgue measurability.\n In fact, together with Borel measurability these are the most prominent regularity conditions studied in descriptive combinatorics.\n \n The corresponding complexity classes of local problems that always admit a solution with the respective regularity property are $\\mathsf{BOREL},\\mathsf{MEASURE},\\mathsf{BAIRE}$ (See \\cref{fig:big_picture_trees}). \n In this paper, we moreover consider the setting where each connected component of the underlying graph is a $\\Delta$-regular tree. \n \n \n \n \n \n \n\t\n\t\n\t\t\n\t\n\t\t\n\t\t\n\n\t The connection between distributed computing and descriptive combinatorics arises from the fact that in descriptive combinatorics we care about constructions that do not use the axiom of choice. \n\t In the distributed language, the axiom of choice corresponds to leader election, that is, the constructions in descriptive combinatorics do not allow picking exactly one point in every component. \n\t \n\t To get some intuition about the power of the complexity class $\\mathsf{BOREL}$, we note that Borel constructions allow us to alternate countably many times the following two operations. \n\t First, any local algorithm with constant local complexity can be run. \n\t Second, we have an oracle that provides a maximal independent set (MIS) on any graph that can be constructed locally from the information computed so far \\cite{KST}. \n Note that from the speedup result of \\cite{chang_kopelowitz_pettie2019exp_separation} we get that every local problem with local complexity $O(\\log^* n)$ can be solved by constant local constructions and \\emph{one} call to such an MIS oracle. This implies the inclusion $\\mathsf{LOCAL}(O(\\log^* n)) \\subseteq \\mathsf{BOREL}$ in \\cref{fig:big_picture_trees} proven in the insightful paper of Bernshteyn \\cite{Bernshteyn2021LLL}. \n The relationship of the class $\\mathsf{MEASURE}$ (and of the class $\\mathsf{fiid}$ from the discussion of factors) to the class $\\mathsf{BOREL}$ is analogous to the relationship of randomized distributed algorithms to deterministic distributed algorithms. \n\n \n \n\n\n\n\n\n\n\t\n\n\t\n\t\n\t\n\n\n\n\n\n\n\\paragraph{Local Problems on Regular Trees}\n\nAfter introducing the three areas of interest in this work, we conclude the section by briefly discussing the kinds of problems we focus on, which are local problems on regular trees.\nMore precisely, we study \\emph{locally checkable labeling (LCL)} problems, which are a class of problems, where the correctness of the solution can be checked locally.\nExamples include classical problems from combinatorics such as proper vertex coloring, proper edge coloring, perfect matching, and maximal independent set.\nOne main goal of this paper is to understand possible complexity classes of LCLs without inputs on infinite $\\Delta$-regular trees and their finite analogues.\nWe refer the reader to \\cref{sec:preliminaries} for a precise definition of a finite $\\Delta$-regular tree.\n\nThe motivation for studying regular trees in this work stems from different sources: (a) infinite $\\Delta$-regular trees are studied in the area of ergodic theory \\cite{Bowen1,Bowen2}, random processes \\cite{backhausz,backhausz2,lyons2011perfect} and descriptive combinatorics \\cite{DetMarks,BrooksMeas}, (b) many lower bounds in distributed computing are proven in regular trees \\cite{balliu_et_al:binary_lcl,balliu2019LB,brandt_etal2016LLL,brandt19automatic,Brandt20tightLBmatching,ChangHLPU20,goeoes14_PODC}, and (c) connecting and comparing the techniques of the three areas in this simple setting reveals already deep connections, see \\cref{sec:contribution}.\n\n\n\n\n\n\n\n\n\n\\section{Our Contributions}\n\\label{sec:contribution}\n\nWe believe that our main contribution is presenting all three perspectives as part of a common theory. Our technical contribution is split into three main parts. \n\n\n\n\n\\subsection{Generalization of Marks' Technique}\\label{sec:intromarks}\nIn Section \\ref{sec:marks} we extend the Borel determinacy technique of Marks \\cite{DetMarks}, which was used to prove the nonexistence of Borel $\\Delta$-colorings and perfect matchings, to a broader class of problems, and \nadapt the extended technique to the distributed setting, thereby obtaining a simple method for proving distributed lower bounds.\nThis method is the first lower bound technique for distributed computing using ideas coming from descriptive combinatorics (see \\cite{BernshteynVizing} for a distributed computing upper bound motivated by descriptive combinatorics).\nMoreover, we show how to use the developed techniques to obtain both $\\mathsf{BOREL}$ and $\\mathsf{LOCAL}$ lower bounds for local problems from a natural class, called homomorphism problems.\nOur key technical ingredient for obtaining the mentioned techniques and results is a novel technique based on the notion of an \\emph{ID graph}.\nWe note that a very similar concept to the ID graph was independently discovered by \\cite{id_graph}. \n\n\\paragraph{Marks' technique}\nIn the following we give an introduction to Marks' technique by going through a variant of his proof~\\cite{DetMarks,Marks_Coloring} that shows that $\\Delta$-coloring has deterministic local complexity $\\Omega(\\log n)$. The proof already works in the case where the considered regular tree comes with an input $\\Delta$-edge coloring. In this case, the output color of a given vertex $u$ can be interpreted as that $u$ ``grabs'' the incident edge of that color. The problem is hence equivalent to the \\emph{edge grabbing} problem where every vertex is required to grab an incident edge such that no two vertices grab the same edge. \n\nWe first show the lower bound in the case that vertices do not have unique identifiers but instead are properly colored with $L > \\Delta$ colors. \nSuppose there is an algorithm $\\mathcal{A}$ solving the edge grabbing problem with local complexity $t(n) = o(\\log n)$, and consider a tree rooted at vertex $u$ of depth $t(n)$; such a tree has less than $n$ vertices, for large enough $n$.\nAssume that $u$ has input color $\\sigma \\in [L]$, and consider, for some fixed edge color $\\alpha$, the edge $e$ that is incident to $u$ and has color $\\alpha$.\nTwo players, Alice and Bob, are playing the following game. In the $i$-th round, Alice colors the vertices at distance $i$ from $u$ in the subtree reachable via edge $e$ with colors from $[L]$. Then, Bob colors all other vertices at distance $i$ from $u$ with colors from $[L]$ (see \\cref{fig:game}).\nConsider the output of $u$ when executing $\\mathcal{A}$ on the obtained colored tree.\nBob wins the game if $u$ grabs the edge $e$, and Alice wins otherwise. \n\nNote that either Alice or Bob has a winning strategy. \nGiven the color $\\sigma$ of $u$, if, for each edge color $\\alpha$, Alice has a winning strategy in the game corresponding to the pair $(\\sigma, \\alpha)$, then we can create $\\Delta$ copies of Alice and let them play their strategy on each subtree of $u$, telling them that the colors chosen by the other Alices are what Bob played. The result is a coloring of the input tree such that $u$, by definition, does not pick any edge, contradicting the fact that $\\mathcal{A}$ provides a valid solution!\nSo for every $\\sigma$ there is at least one $\\alpha$ such that Bob has a winning strategy for the game corresponding to $(\\sigma, \\alpha)$. By the pigeonhole principle, there are two colors $\\sigma_1, \\sigma_2$, such that Bob has a winning strategy for both pairs $(\\sigma_1, \\alpha)$ and $(\\sigma_2, \\alpha)$. \nBut now we can imagine a tree rooted in an edge between vertices $u_1, u_2$ that are colored with colors $\\sigma_1, \\sigma_2$. We can now take two copies of Bob, one playing at $u_1$ and the other playing at $u_2$ and let them battle it out, telling each copy that the other color from $\\{\\sigma_1, \\sigma_2\\}$ and whatever the other copy plays are the moves of Alice. The resulting coloring has the property that both $u_1$ and $u_2$, when executing $\\mathcal{A}$ on the obtained colored tree, grab the edge between them, a contradiction that finishes the proof!\n\n\\paragraph{The ID graph}\nThe downside of the proof is that it does not work in the model with unique identifiers (where the players' moves consist in assigning identifiers instead of colors), since gluing copies of the same player could result in an identifier assignment where the identifiers are not unique. \nOne possible remedy is to conduct the whole proof in the context of Borel graphs as was done by Marks. This proves an even stronger statement, namely that $\\Delta$-coloring is not in the class $\\mathsf{BOREL}$, but requires additional ad-hoc tricks and a pretty heavy set theoretic tool---Martin's celebrated Borel determinacy theorem~\\cite{martin} stating that even for infinite two-player games one of the players has to have a winning strategy if the payoff set is Borel. \nThe ID graph enables us to adapt the proof (and its generalization that we develop in \\cref{sec:marks}) to the distributed setting, where the fact that one of the players has a winning strategy is obvious. \nMoreover, we use an infinite version of the ID graph to generalize Marks' technique also in the Borel setting. \n\nHere is how it works: \nThe ID graph is a specific graph whose vertices are the possible unique input identifiers (for input graphs of size $n$), that is, numbers from $[n^{O(1)}]$. Its edges are colored with colors from $[\\Delta]$ and its girth is $\\Omega(\\log n)$. When we define the game between Alice and Bob, we require them to label vertices with identifiers in such a way that whenever a new vertex is labeled with identifier $i$, and its already labeled neighbor has identifier $j$, then $ij$ is an edge in the ID graph.\nMoreover, the color of edge $ij$ in the ID graph is required to be the same as the color of the edge between the vertices labeled $i$ and $j$ in the tree where the game is played. \nIt is straightforward to check that these conditions enforce that even if we let several copies of the same player play, the resulting tree is labeled with unique identifiers. Hence, the same argument as above now finally proves that the deterministic local complexity of $\\Delta$-coloring is $\\Omega (\\log n)$. \n\nWe note that our ID graph technique is of independent interest and may have applications in many different contexts. \nTo give an example from distributed computing, consider the proof of the deterministic $\\Omega(\\log n)$-round lower bound for $\\Delta$-coloring developed by the distributed community~\\cite{brandt_etal2016LLL,chang_kopelowitz_pettie2019exp_separation}, which is based on the celebrated round elimination technique. \nEven though the result is deterministic, the proof is quite technical due to the fact that it relies on examining randomized algorithms, for reasons similar to the reasons why Marks' proof does not apply directly to the setting with unique identifiers. \nFortunately, it can be again streamlined with the use of the ID graph technique. \nMoreover, in a follow-up work, the ID graph technique has already led to a new lower bound for the Lov\u00e1sz local lemma~\\cite{brandt_grunau_rozhon2021LLL_in_LCA} in the area of Local Computation Algorithms (which is part of the realm of sequential computation), thereby giving further evidence for the versatility of the technique. \n\n\n\\paragraph{Marks vs. Round Elimination}\nIt is quite insightful to compare Marks' technique (and our generalization of it) with the powerful round elimination technique~\\cite{brandt19automatic}, which has been responsible for all locality lower bounds of the last years~\\cite{brandt_etal2016LLL, balliu2019LB, brandt19automatic, balliu2019hardness_homogeneous, balliu_et_al:binary_lcl, Brandt20tightLBmatching, balliu20rulingset, ChangHLPU20, balliu21dominatingset}. \nWhile, on the surface, Marks' approach developed for the Borel world may seem quite different from the round elimination technique, there are actually striking similarities between the two methods.\nOn a high level, in the round elimination technique, the following argument is used to prove lower bounds in the LOCAL model:\nIf a $T$-round algorithm exists for a problem $\\Pi_0$ of interest, then there exists a $(T-1)$-round algorithm for some problem $\\Pi_1$ that can be obtained from $\\Pi_0$ in a mechanical manner.\nBy applying this step iteratively, we obtain a problem $\\Pi_t$ that can be solved in $0$ rounds; by showing that there is no $0$-algorithm for $\\Pi_t$ (which is easy to do if $\\Pi_t$ is known), a $(T+1)$-round lower bound for $\\Pi_0$ is obtained.\n\nThe interesting part regarding the relation to Marks' technique is how the ($T-i-1$)-round algorithms $\\mathcal{A}'$ are obtained from the ($T-i$)-round algorithms $\\mathcal{A}$ in the round elimination framework: in order to execute $\\mathcal{A}'$, each vertex $v$, being aware of its ($T-i-1$)-hop neighborhood, essentially asks whether, for all possible extensions of its view by one hop along a chosen incident edge, there exists some extension of its view by one hop along all other incident edges such that $\\mathcal{A}$, executed on the obtained ($T-i$)-hop neighborhood, returns a certain output at $v$, and then bases its output on the obtained answer.\nIt turns out that the vertex sets corresponding to these two extensions correspond precisely to two moves of the two players in the game(s) played in Marks' approach: more precisely, in round $T-i$ of a game corresponding to the considered vertex $v$ and the chosen incident edge, the move of Alice consists in labeling the vertices corresponding to the first extension, and the move of Bob consists in labeling the vertices corresponding to the second extension.\n\nHowever, despite the similarities, the two techniques (at least in their current forms) have their own strengths and weaknesses and are interestingly different in that there are local problems that we know how to obtain good lower bounds for with one technique but not the other, and vice versa.\nFinding provable connections between the two techniques is an exciting research direction that we expect to lead to a better understanding of the possibilities and limitations of both techniques.\n\nIn \\cref{sec:marks} we use our generalized and adapted version of Marks' technique to prove new lower bounds for so-called homomorphism problems. \nHomomorphism problems are a class of local problems that generalizes coloring problems---each vertex is to be colored with some color and there are constraints on which colors are allowed to be adjacent. \nThe constraints can be viewed as a graph---in the case of coloring this graph is a clique. \nIn general, whenever the underlying graph of the homomorphism problem is $\\Delta$-colorable, its deterministic local complexity is $\\Omega(\\log n)$, because solving the problem would imply that we can solve $\\Delta$-coloring too (in the same runtime). \nIt seems plausible that homomorphism problems of this kind are the only hard, i.e., $\\Omega(\\log n)$, homomorphism problems. However, our generalization of Marks' technique asserts that this is not true. \n\n\\begin{theorem}\nThere are homomorphism problems whose deterministic local complexity on trees is $\\Omega(\\log n)$ such that the chromatic number of the underlying graph is $2\\Delta-2$.\n\\end{theorem}\n\nIt is not known how to prove the same lower bounds using round elimination\\footnote{Indeed, the descriptions of the problems have comparably large numbers of labels and do not behave like so-called ``fixed points'' (i.e., nicely) under round elimination, which suggests that it is hard to find a round elimination proof with the currently known approaches.}; in fact, as far as we know, these problems are the only known examples of problems on $\\Delta$-regular trees for which a lower bound is known to hold but currently not achievable by round elimination. \nProving the same lower bounds via round elimination is an exciting open problem. \n\n\n\n\n\n\n\n\n\n\n\\subsection{Separation of Various Complexity Classes}\n\\label{sec:introseparation}\n\\paragraph{Uniform Complexity Landscape}\nWe investigate the connection between randomized and uniform distributed local algorithms, where uniform algorithms are equivalent to the studied notion of finitary factors of iid. \nFirst, it is simple to observe that local problems with uniform complexity $t(\\varepsilon)$ have randomized complexity $t(1 \/ n^{O(1)})$ -- by definition, every vertex knows its local output after that many rounds with probability $1 - 1\/n^{O(1)}$. The result thus follows by a union bound over the $n$ vertices of the input graph.\n\nOn the other hand, we observe that on $\\Delta$-regular trees the implication also goes in the opposite direction in the following sense. Every problem that has a randomized complexity of $t(n) = o(\\log n)$ has an uniform complexity of $O(t(1\/\\varepsilon))$. \n\nOne could naively assume that this equivalence also holds for higher complexities, but this is not the case. Consider for example the $3$-coloring problem.\nIt is well-known in the distributed community that $3$-coloring a tree can be solved deterministically in $O(\\log n)$ rounds using the rake-and-compress decomposition~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup,MillerR89}.\nOn the other hand, there is no uniform algorithm for $3$-coloring a tree. \nIf there were such an uniform algorithm, we could run it on any graph with large enough girth and color $99\\%$ of its vertices with three colors. This in turn would imply that the high-girth graph has an independent set of size at least $0.99 \\cdot n\/3$. This is a contradiction with the fact that there exists high-girth graphs with a much smaller independence number~\\cite{bollobas}. \n\nInterestingly, the characterization of Bernshteyn \\cite{Bernshteyn_work_in_progress} implies that \\emph{any} uniform distributed algorithm can be ``sped up'' to a deterministic local $O(\\log n)$ complexity, as we prove in~\\cref{thm:MainBaireLog}. \n\nWe show that there are local problems that can be solved by an uniform algorithm but only with a complexity of $\\Omega(\\log 1\/\\varepsilon)$. Namely, the problem of constructing a $2$-hop perfect matching on infinite $\\Delta$-regular trees for $\\Delta \\ge 3$ has an uniform local complexity between $\\Omega(\\log 1\/\\varepsilon)$ and $O(\\operatorname{poly}\\log 1\/\\varepsilon)$. \nFormally, this proves the following theorem. \n\n\\begin{theorem}\n$\\mathsf{OLOCAL}(O(\\log\\log 1\/\\varepsilon)) \\subsetneq \\mathsf{OLOCAL}(O(\\operatorname{poly}\\log 1\/\\varepsilon))$. \n\\end{theorem}\n\nThe uniform algorithm for this problem is based on a so-called one-ended forest decomposition introduced in \\cite{BrooksMeas} in the descriptive combinatorics context. In a one-ended forest decomposition, each vertex selects exactly one of its neighbors as its parent by orienting the corresponding edge outwards. This defines a decomposition of the vertices into infinite trees. We refer to such a decomposition as a one-ended forest decomposition if the subtree rooted at each vertex only contains finitely many vertices. \nHaving computed such a decomposition, $2$-hop perfect matching can be solved inductively starting from the leaf vertices of each tree.\n\nWe leave the understanding of the uniform complexity landscape in the regime $\\Omega(\\log 1\/\\varepsilon)$ as an exciting open problem. In particular, does there exist a function $g(\\varepsilon)$ such that each local problem that can be solved by an uniform algorithm has an uniform complexity of $O(g(\\varepsilon))$?\n\n\n\n\\paragraph{Relationship of Distributed Classes with Descriptive Combinatorics}\n\nBernshteyn recently proved that $\\mathsf{LOCAL}(\\log^* n) \\subseteq \\mathsf{BOREL}$ \\cite{Bernshteyn2021LLL}. That is, each local problem with a deterministic $\\mathsf{LOCAL}$ complexity of $O(\\log^* n)$ also admits a Borel-measurable solution. A natural question to ask is whether the converse also holds. \nIndeed, it is known that $\\mathsf{LOCAL}(\\log^* n) = \\mathsf{BOREL}$ on paths with no additional input \\cite{grebik_rozhon2021LCL_on_paths}. \nWe show that on regular trees the situation is different. \nOn one hand, a characterization of Bernshteyn \\cite{Bernshteyn_work_in_progress} implies that $\\mathsf{BOREL} \\subseteq \\mathsf{BAIRE} \\subseteq \\mathsf{LOCAL}(O(\\log n))$. \nOn the other hand, we show that this result cannot be strengthened by proving the following result. \n\n\\begin{theorem}\n$\\mathsf{BOREL} \\subsetneq \\mathsf{RLOCAL}(o(\\log n))$.\n\\end{theorem}\n\nThat is, there exists a local problem that admits a Borel-measurable solution but cannot be solved with a (randomized) $\\mathsf{LOCAL}$ algorithm running in a sublogarithmic number of rounds. \n\nLet us sketch a weaker separation, namely that $\\mathsf{BOREL} \\setminus \\mathsf{LOCAL}(O(\\log^* n)) \\neq \\emptyset$. Consider a version of $\\Delta$-coloring where a subset of vertices can be left uncolored. However, the subgraph induced by the uncolored vertices needs to be a collection of doubly-infinite paths (in finite trees, this means each path needs to end in a leaf vertex). \nThe nonexistence of a fast distributed algorithm for this problem essentially follows from the celebrated $\\Omega(\\log n)$ deterministic lower bound for $\\Delta$-coloring of~\\cite{brandt_etal2016LLL}. \nOn the other hand, the problem allows a Borel solution. First, sequentially compute $\\Delta-2$ maximal independent sets, each time coloring all vertices in the MIS with the same color, followed by removing all the\ncolored vertices from the graph. In that way, a total of $\\Delta-2$ colors are used. Moreover, each\nuncolored vertex has at most $2$ uncolored neighbors. \nThis implies that the set of uncolored vertices forms a disjoint union of finite paths, one ended\ninfinite paths and doubly infinite paths. The first two classes can be colored inductively with two\nadditional colors, starting at one endpoint of each path in a Borel way (namely it can be done by making use of the countably many MISes in larger and larger powers of the input graph). \nHence, in the end only doubly infinite paths are left uncolored, as desired. \n\nTo show the stronger separation between the classes $\\mathsf{BOREL}$ and $\\mathsf{RLOCAL}(o(\\log n))$ we use a variation of the\n$2$-hop perfect matching problem. In this variation, some of the vertices can be left unmatched, but similar as in the variant of the $\\Delta$-coloring problem described above, the graph induced by all the unmatched vertices needs to satisfy some additional constraints.\n\nWe conclude the paragraph by noting that the separation between the classes $\\mathsf{BOREL}$ and $\\mathsf{LOCAL}(O(\\log^* n))$ is not as simple as it may look in the following sense. \nThis is because problems typically studied in the $\\mathsf{LOCAL}$ model with a $\\mathsf{LOCAL}$ complexity of $\\omega(\\log^* n)$ like $\\Delta$-coloring and perfect matching also do not admit a Borel-measurable solution due to the technique of Marks \\cite{DetMarks} that we discussed in \\cref{sec:intromarks}.\n\n\n\n\n\n\\subsection{$\\mathsf{LOCAL}(O(\\log n)) = \\mathsf{BAIRE}$} \n\nWe already discussed that one of complexity classes studied in descriptive combinatorics is the class $\\mathsf{BAIRE}$.\nRecently, Bernshteyn proved \\cite{Bernshteyn_work_in_progress} that all local problems that are in the complexity class $\\mathsf{BAIRE}, \\mathsf{MEASURE}$ or $\\mathsf{fiid}$ have to satisfy a simple combinatorial condition which we call being $\\ell$-full. \nOn the other hand, all $\\ell$-full problems allow a $\\mathsf{BAIRE}$ solution \\cite{Bernshteyn_work_in_progress}.\nThis implies a complete combinatorial characterization of the class $\\mathsf{BAIRE}$. \nWe defer the formal definition of $\\ell$-fullness to Section \\ref{sec:baire} as it requires a formal definition of a local problem. \nInformally speaking, in the context of vertex labeling problems, a problem is $\\ell$-full if we can choose a subset $S$ of the labels with the following property.\nWhenever we label two endpoints of a path of at least $\\ell$ vertices with two labels from $S$, we can extend the labeling with labels from $S$ to the whole path such that the overall labeling is valid.\nFor example, proper $3$-coloring is $3$-full with $S=\\{1,2,3\\}$ because for any path of three vertices such that its both endpoints are colored arbitrarily, we can color the middle vertex so that the overall coloring is proper. \nOn the other hand, proper $2$-coloring is not $\\ell$-full for any $\\ell$. \n\nWe complement this result as follows. \nFirst, we prove that any $\\ell$-full problem has local complexity $O(\\log n)$, thus proving that all complexity classes considered in the areas of factors of iids and descriptive combinatorics from \\cref{fig:big_picture_trees} are contained in $\\mathsf{LOCAL}(O(\\log n))$. \nIn particular, this implies that the existence of \\emph{any} uniform algorithm implies a local distributed algorithm for the same problem of local complexity $O(\\log n)$. \nWe obtain this result via the well-known rake-and-compress decomposition~\\cite{MillerR89}. \n\n\nOn the other hand, we prove that any problem in the class $\\mathsf{LOCAL}(O(\\log n))$ satisfies the $\\ell$-full condition. \nThe proof combines a machinery developed by Chang and Pettie~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}\nwith additional nontrivial ideas.\nIn this proof we construct recursively a sequence of sets of rooted, layered, and partially labeled trees, where the partial labeling is computed by simulating any given $O(\\log n)$-round distributed algorithm, \nand then the set $S$ meeting the $\\ell$-full condition is constructed by considering all possible extensions of the partial labeling to complete correct labeling of these trees.\n\nThis result implies the following equality:\n\\begin{theorem}\n$\\mathsf{LOCAL}(O(\\log n)) = \\mathsf{BAIRE}$. \n\\end{theorem}\n\nThis equality is surprising in that the definitions of the two classes do not seem to have much in common on the first glance!\nMoreover, the proof of the equality relies on nontrivial results in both distributed algorithms (the technique of Chang and Pettie~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}) and descriptive combinatorics (the fact that a hierarchical decomposition, so-called toast, can be constructed in $\\mathsf{BAIRE}$, \\cite{conleymillerbound}, see Proposition \\ref{pr:BaireToast}). \n\n\n\n\n\n\n\n\n\n\n\n\nThe combinatorial characterization of the local complexity class $\\mathsf{LOCAL}(O(\\log n))$ on $\\Delta$-regular trees is interesting from the perspective of distributed computing alone. \nThis result can be seen as a part of a large research program aiming at classification of possible local complexities on various graph classes~\\cite{balliu2020almost_global_problems,brandt_etal2016LLL,chang_kopelowitz_pettie2019exp_separation,chang_pettie2019time_hierarchy_trees_rand_speedup,chang2020n1k_speedups,Chang_paths_cycles,balliu2021_rooted_trees,balliu_et_al:binary_lcl,balliu2018new_classes-loglog*-log*}. \nThat is, we wish not only to understand possible complexity classes (see the left part of \\cref{fig:big_picture_trees} for possible local complexity classes on regular trees), but also to find combinatorial characterizations of problems in those classes that allow us to efficiently decide for a given problem which class it belongs to. \nUnfortunately, even for grids with input labels, it is \\emph{undecidable} whether a given local problem can be solved in $O(1)$ rounds~\\cite{naorstockmeyer,brandt_grids}, since\nlocal problems on grids can be used to simulate a Turing machine.\nThis undecidability result does not apply to paths and trees, hence for these graph classes it is still hopeful that we can find simple and useful characterizations for different classes of distributed problems. \n\nIn particular, on paths it is decidable what classes a given local problem belongs to, for all classes coming from the three areas considered here, and this holds even if we allow inputs~\\cite{Chang_paths_cycles,grebik_rozhon2021LCL_on_paths}. \nThe situation becomes much more complicated when we consider trees. \nRecently, classification results on trees were obtained for so-called binary-labeling problems \\cite{balliu_et_al:binary_lcl}. \nMore recently, a complete classification was obtained in the case of \\emph{rooted regular trees}~\\cite{balliu2021_rooted_trees}. \nAlthough their algorithm takes exponential time in the worst case, the authors provided a practical implementation fast enough to classify many typical problems of interest. \n\nMuch less is known for general, \\emph{unoriented} trees, with an arbitrary number of labels. \nIn general, deciding the optimal distributed complexity for a local problem on bounded-degree trees is $\\mathsf{EXPTIME}$-hard~\\cite{chang2020n1k_speedups}, such a hardness result does not rule out the possibility for having a simple and polynomial-time characterization for the case of \\emph{regular trees}, where there is no input and the constraints are placed only on degree-$\\Delta$ vertices. \nIndeed, it was stated in~\\cite{balliu2021_rooted_trees} as an open question to find such a characterization. Our characterization of $\\mathsf{LOCAL}(O(\\log n)) = \\mathsf{BAIRE}$ by $\\ell$-full problems makes progress in better understanding the distributed complexity classes on trees and towards answering this open question. \n\n\n\n\n\n\n\n\\paragraph{Roadmap}\n\\hypertarget{par:roadmap}{}\nIn \\cref{sec:preliminaries}, we define formally all the three setups we consider in the paper. \nIn \\cref{sec:marks} we discuss the lower bound technique of Marks and the new concept of an ID graph.\nNext, in \\cref{sec:separating_examples} we prove some basic results about the uniform complexity classes and give examples of problems separating some classes from \\cref{fig:big_picture_trees}. \nFinally, in \\cref{sec:baire} we prove that a problem admits a Baire measurable solution if and only if it admits a distributed algorithm of local complexity $O(\\log n)$.\n\nThe individual sections can be read largely independently of each other.\nMoreover, most of our results that are specific to only one of the three areas can be understood without reading the parts of the paper that concern the other areas.\nWe encourage the reader interested mainly in one of the areas to skip the respective parts.\n\n\n\n\n\\section{Preliminaries}\n\\label{sec:preliminaries}\n\nIn this section, we explain the setup we work with, the main definitions and results. \nThe class of graphs that we consider in this work are either infinite $\\Delta$-regular trees, or their finite analogue that we define formally in \\cref{subsec:local_problems}.\n\nWe sometimes explicitly assume $\\Delta > 2$. \nThe case $\\Delta = 2$, that is, studying paths, behaves differently and seems much easier to understand \\cite{grebik_rozhon2021LCL_on_paths}. \nUnless stated otherwise, we do not consider any additional structure on the graphs, but sometimes it is natural to work with trees with an input $\\Delta$-edge-coloring. \n\n\n\\subsection{Local Problems on $\\Delta$-regular trees}\n\\label{subsec:local_problems}\n\nThe problems we study in this work are locally checkable labeling (LCL) problems, which, roughly speaking, are problems that can be described via local constraints that have to be satisfied in a suitable neighborhood of each vertex.\nIn the context of distributed algorithms, these problems were introduced in the seminal work by Naor and Stockmeyer~\\cite{naorstockmeyer}, and have been studied extensively since.\nIn the modern formulation introduced in~\\cite{brandt19automatic}, instead of labeling vertices or edges, LCL problems are described by labeling half-edges, i.e., pairs of a vertex and an incident edge.\nThis formulation is very general in that it not only captures vertex and edge labeling problems, but also others such as orientation problems, or combinations of all of these types.\nBefore we can provide this general definition of an LCL, we need to introduce some definitions.\nWe start by formalizing the notion of a half-edge.\n\n\\begin{definition}[Half-edge]\n\\label{def:half-edge}\n A \\emph{half-edge} is a pair $(v, e)$ where $v$ is a vertex, and $e$ an edge incident to $v$.\n We say that a half-edge $(v, e)$ is \\emph{incident} to a vertex $w$ if $w = v$, we say that a vertex $w$ is \\emph{contained} in a half edge $(v,e)$ if $w=v$, and we say that $(v, e)$ \\emph{belongs} to an edge $e'$ if $e' = e$.\n We denote the set of all half-edges of a graph $G$ by $H(G)$.\n A \\emph{half-edge labeling} is a function $c \\colon H(G) \\to \\Sigma$ that assigns to each half-edge an element from some label set $\\Sigma$.\n\\end{definition}\n\nIn order to speak about finite $\\Delta$-regular trees, we need to consider slightly modified definition of a graph.\nWe think of each vertex to be contained in $\\Delta$-many half-edges, however, not every half edge belongs to an actual edge of the graph.\nThat is half-edges are pairs $(v,e)$, but $e$ is formally not a pair of vertices.\nSometimes we refer to these half-edges as \\emph{virtual} half-edges.\nWe include a formal definition to avoid confusions. See also \\cref{fig:Delta_reg_tree}. \n\n\n\\begin{definition}[$\\Delta$-regular trees]\n A tree $T$, finite or infinite is a \\emph{$\\Delta$-regular tree} if either it is infinite and $T=T_\\Delta$, where $T_\\Delta$ is the unique infinite $\\Delta$-regular tree, that is each vertex has exactly $\\Delta$-many neighbors, or it is finite of maximum degree $\\Delta$ and each vertex $v\\in T$ of degree $d\\le \\Delta$ is contained in $(\\Delta-d)$-many virtual half-edges.\n \n \n Formally, we can view $T$ as a triplet $(V(T), E(T), H(T))$, where $(V(T),E(T))$ is a tree of maximum degree $\\Delta$ and $H(T)$ consists of real half-edges, that is pairs $(v,e)$, where $v\\in V(T)$, $e\\in E(T)$ and $e$ is incident to $v$, together with some virtual edges, in the case when $T$ is finite, such that each vertex $v\\in V(T)$ is contained in exactly $\\Delta$-many half-edges (real or virtual).\n \n \n\\end{definition}\n\n\\begin{figure}\n \\centering\n \\includegraphics{fig\/Delta_reg_tree.pdf}\n \\caption{A $3$-regular tree on 15 vertices. Every vertex has $3$ half-edges but not all half-edges lead to a different vertex. }\n \\label{fig:Delta_reg_tree}\n\\end{figure}\n\nAs we are considering trees in this work, each LCL problem can be described in a specific form that provides two lists, one describing all label combinations that are allowed on the half-edges incident to a vertex, and the other describing all label combinations that are allowed on the two half-edges belonging to an edge.\\footnote{Every problem that can be described in the form given by Naor and Stockmeyer~\\cite{naorstockmeyer} can be equivalently described as an LCL problem in this list form, by simply requiring each output label on some half-edge $h$ to encode all output labels in a suitably large (constant) neighborhood of $h$ in the form given in~\\cite{naorstockmeyer}.}\nWe arrive at the following definition for LCLs on $\\Delta$-regular trees.\\footnote{Note that the defined LCL problems do not allow the correctness of the output to depend on input labels.}\n\n\\begin{definition}[LCLs on $\\Delta$-regular trees]\\label{def:LCL_trees}\nA \\emph{locally checkable labeling problem}, or \\emph{LCL} for short, is a triple $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$, where $\\Sigma$ is a finite set of labels, $\\mathcal{V}$ is a subset of unordered cardinality-$\\Delta$ multisets\\footnote{Recall that a multiset is a modification of the concept of sets, where repetition is allowed.} of labels from $\\Sigma$ , and $\\mathcal{E}$ is a subset of unordered cardinality-$2$ multisets of labels from $\\Sigma$.\n\nWe call $\\mathcal{V}$ and $\\mathcal{E}$ the \\emph{vertex constraint} and \\emph{edge constraint} of $\\Pi$, respectively.\nMoreover, we call each multiset contained in $\\mathcal{V}$ a \\emph{vertex configuration} of $\\Pi$, and each multiset contained in $\\mathcal{E}$ an \\emph{edge configuration} of $\\Pi$.\n\nLet $T$ be a $\\Delta$-regular tree and $c:H(T)\\to \\Sigma$ a half-edge labeling of $T$ with labels from $\\Sigma$.\nWe say that $c$ is a \\emph{$\\Pi$-coloring}, or, equivalently, a correct solution for $\\Pi$, if, for each vertex $v$ of $T$, the multiset of labels assigned to the half-edges incident to $v$ is contained in $\\mathcal{V}$, and, for each edge $e$ of $T$, the cardinality-$2$ multiset of labels assigned to the half-edges belonging to $e$ is an element of $\\mathcal{E}$.\n\\end{definition}\n\nAn equivalent way to define our setting would be to consider $\\Delta$-regular trees as commonly defined, that is, there are vertices of degree $\\Delta$ and vertices of degree $1$, i.e., leaves.\nIn the corresponding definition of LCL one would consider leaves as unconstrained w.r.t. the vertex constraint, i.e., in the above definition of a correct solution the condition ``for each vertex $v$'' is replaced by ``for each non-leaf vertex $v$''.\nEquivalently, we could allow arbitrary trees of maximum degree $\\Delta$ as input graphs, but, for vertices of degree $< \\Delta$, we require the multiset of labels assigned to the half-edges to be extendable to some cardinality-$\\Delta$ multiset in $\\mathcal{V}$.\nWhen it helps the exposition of our ideas and is clear from the context, we may make use of these different but equivalent perspectives.\n\nWe illustrate the difference between our setting and ``standard setting'' without virtual half-edges on the perfect matching problem.\nA standard definition of the perfect matching problem is that some edges are picked in such a way that each vertex is covered by exactly one edge.\nIt is easy to see that there is no local algorithm to solve this problem on the class of finite trees (without virtual half-edges), this is a simple parity argument.\nHowever, in our setting, every vertex needs to pick exactly one half-edge (real or virtual) in such a way that both endpoints of each edge are either picked or not picked.\nWe remark that in our setting it is not difficult to see that (if $\\Delta>2$), then this problem can be solved by a local deterministic algorithm of local complexity $O(\\log(n))$.\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{The $\\mathsf{LOCAL}$ model}\n\nIn this section, we define local algorithms and local complexity.\nWe discuss the general relation between the classical local complexity and the uniform local complexity. Recall that when we talk about distributed algorithm on $\\Delta$-regular trees, the algorithm has access to $n$, the size of the tree. The measure of complexity is the classical local complexity. On the other hand, when we talk about uniform distributed algorithms on $\\Delta$-regular trees, we talk about an infinite $\\Delta$-regular tree and the measure of complexity is the uniform local complexity. \nWe start with the classical notions.\nRecall that $B(v,t)$ is the $t$-hop neighborhood of a vertex $v$ in an underlying graph $G$.\n\n\\begin{definition}[Local algorithm]\n\\label{def:local_algorithm}\nA \\emph{distributed local algorithm} $\\mathcal{A}$ of local complexity $t(n)$ is a function defined on all possible $t(n)$-hop neighborhoods of a vertex. Applying an algorithm $\\mathcal{A}$ on an input graph $G$ means that the function is applied to a $t(n)$-hop neighborhood $B(u, t(n))$ of each vertex $u$ of $G$. The output of the function is a labeling of the half-edges around the given vertex. The algorithm also takes as input the size of the input graph $n$. \n\\end{definition}\n\n\\begin{definition}[Local complexity]\n\\label{def:local_complexity}\nWe say that an LCL problem $\\Pi$ has a deterministic local complexity $t(n)$ if there is a local algorithm $\\mathcal{A}$ of local complexity $t(n)$ such that when run on the input graph $G$, with each of its vertices having a unique identifier from $[n^{O(1)}]$, $\\mathcal{A}$ always returns a valid solution to $\\Pi$. \nWe also say $\\Pi \\in \\mathsf{DLOCAL}( O(t(n))$. \n\nThe randomized complexity is define analogously, but instead of unique identifiers, each vertex of $G$ has an infinite random string. The solution of $\\mathcal{A}$ needs to be a valid solution with probability $1 - 1\/n$. \nWe also say $\\Pi \\in \\mathsf{RLOCAL}( O(t(n))$. \n\\end{definition}\n\nWhenever we talk about \\emph{local} complexity on $\\Delta$-regular trees, we always tacitly think about the class of finite $\\Delta$-regular trees. We use the notation $\\mathsf{LOCAL}(O(t(n)))$ whenever the deterministic and the randomized complexity of the problems are the same up to constant factor (cf. \\cref{thm:basicLOCAL}). \n\n\\paragraph{Classification of Local Problems on $\\Delta$-regular Trees}\n\nThere is a lot of work aiming to classify all possible local complexities of LCL problems on bounded degree trees. This classification was recently finished.\nEven though the results in the literature are stated for the classical notion of finite trees, we note that it is easy to check that the same picture emerges if we restrict ourselves to $\\Delta$-regular trees according to our definition.\n\n\\begin{theorem}[Classification of local complexities of LCLs on $\\Delta$-regular trees \\cite{naorstockmeyer,chang_kopelowitz_pettie2019exp_separation,chang_pettie2019time_hierarchy_trees_rand_speedup,chang2020n1k_speedups,balliu2020almost_global_problems,balliu2019hardness_homogeneous,brandt_grunau_rozhon2021classification_of_lcls_trees_and_grids}]\n\\label{thm:basicLOCAL}\nLet $\\Pi$ be an LCL problem. Then the deterministic\/randomized local complexity of $\\Pi$ on $\\Delta$-regular trees is one of the following:\n\\begin{enumerate}\n \\item $O(1)$,\n \\item $\\Theta(\\log^*n)$,\n \\item $\\Theta(\\log n)$ deterministic and $\\Theta(\\log\\log n)$ randomized,\n \\item $\\Theta(\\log n)$,\n \\item $\\Theta(n^{1\/k})$ for $k \\in \\mathbb{N}$.\n\\end{enumerate}\n\\end{theorem}\n\nWe do not know whether the complexity classes $\\Theta(n^{1\/k})$ are non-empty on \\emph{$\\Delta$-regular} trees. The $2 \\frac{1}{2}$ coloring problems $\\mathcal{P}_k$ are known to have complexity $\\Theta(n^{1\/k})$ on \\emph{bounded-degree} trees~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}. However, as the correctness criterion of $\\mathcal{P}_k$ depends on the degree of vertices, it does not fit into the class of LCL problems that we consider in \\cref{def:LCL_trees}. \nNevertheless, as can be seen in \\cref{fig:big_picture_trees}, all classes of factors of iid and from descriptive combinatorics are already contained in the class $\\mathsf{LOCAL}(O(\\log n))$. \n\n\n\n\\paragraph{Uniform Algorithms}\nAs we mentioned before, when talking about local complexities, we always have in mind that the underlying graph is finite.\nIn particular, the corresponding algorithm knows the size of the graph.\nOn infinite $\\Delta$-regular trees, or infinite graphs in general, we use the following notion \\cite{HolroydSchrammWilson2017FinitaryColoring,Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring}. \n\n\\begin{definition}\nAn \\emph{uniform local algorithm} $\\mathcal{A}$ is a function that is defined on all possible (finite) neighborhoods of a vertex.\nFor some neighborhoods it outputs a special symbol $\\emptyset$ instead of a labeling of the half-edges around the central vertex.\nApplying $\\mathcal{A}$ on a graph $G$ means that for each vertex $u$ of $G$ the function is applied to $B(u,t)$, where $t$ is the minimal number such that $\\mathcal{A}(u,t)\\not=\\emptyset$.\nWe call $t$ the coding radius of $\\mathcal{A}$, and denote it, as a function on vertices, as $R_\\mathcal{A}$.\n\\end{definition}\n\nWe define the corresponding notion of \\emph{uniform local complexity} for infinite $\\Delta$-regular trees where each vertex is assigned an infinite random string.\n\n\\begin{definition}[Uniform local complexity \\cite{HolroydSchrammWilson2017FinitaryColoring}]\n\\label{def:uniform_complexity}\nWe say that the \\emph{uniform local (randomized) complexity} of an LCL problem $\\Pi$ is $t(\\varepsilon)$ if there is an uniform local algorithm $\\mathcal{A}$ such that the following hold on the infinite $\\Delta$-regular tree.\nRecall that $R_\\mathcal{A}$ is the random variable measuring the coding radius of a vertex $u$, that is, the distance $\\mathcal{A}$ needs to look at to decide the answer for $u$. Then, for any $0 < \\varepsilon < 1$:\n\\[\n\\P(R_\\mathcal{A} \\ge t(\\varepsilon)) \\leq \\varepsilon.\n\\]\nWe also say $\\Pi \\in \\mathsf{OLOCAL}(O(t(\\varepsilon))$.\n\\end{definition}\n\nWe finish by stating the following lemma that bounds the uniform complexity of concatenation of two uniform algorithm (we need to be little bit more careful and cannot just add the complexites up). \n\n\\begin{lemma}[Sequential Composition]\n\\label{lem:sequential_composition}\nLet $A_1$ and $A_2$ be two distributed uniform algorithms with an uniform local complexity of $t_1(\\varepsilon)$ and $t_2(\\varepsilon)$, respectively. \nLet $A$ be the sequential composition of $A_1$ and $A_2$. That is, $A$ needs to know the output that $A_2$ computes when the local input at every vertex is equal to the output of the algorithm $A_1$ at that vertex. Then, $t_A(\\varepsilon) \\leq t_{A_1} \\left( \\frac{\\varepsilon\/2}{\\Delta^{t_{A_2}(\\varepsilon\/2) + 1}}\\right) + t_{A_2}(\\varepsilon\/2)$.\n\\end{lemma}\n\\begin{proof}\nConsider some arbitrary vertex $u$. Let $E_1$ denote the event that the coding radius of $\\mathcal{A}_1$ is at most $t_{A_1} \\left( \\frac{\\varepsilon\/2}{\\Delta^{t_{A_2}(\\varepsilon\/2) + 1}}\\right)$ for all vertices in the $t_{A_2}(\\varepsilon\/2)$-hop neighborhood around $u$. As the $t_{A_2}(\\varepsilon\/2)$-hop neighborhood around $u$ contains at most $\\Delta^{t_{A_2}(\\varepsilon\/2) + 1}$ vertices, a union bound implies that $P(E_1) \\geq 1 - \\varepsilon\/2$. Moreover, let $E_2$ denote the event that the coding radius of algorithm $\\mathcal{A}_2$ at vertex $u$ is at most $t_{A_2}(\\varepsilon\/2)$. By definition, $P(E_2) \\geq 1 - \\varepsilon\/2$. Moreover, if both events $E_1$ and $E_2$ occur, which happens with probability at least $1 - \\varepsilon$, then the coding radius of algorithm $\\mathcal{A}$ is at most $t_{A_1} \\left( \\frac{\\varepsilon\/2}{\\Delta^{t_{A_2}(\\varepsilon\/2) + 1}}\\right) + t_{A_2}(\\varepsilon\/2)$, thus finishing the proof.\n \\end{proof}\n\n\n\\subsection{Descriptive combinatorics}\n\\label{subsec:descriptivecombinatorics}\n\nBefore we define formally the descriptive combinatorics complexity classes, we give a high-level overview on their connection to distributing computing for the readers more familiar with the latter.\n\nThe complexity class that partially captures deterministic local complexity classes is called $\\mathsf{BOREL}$ (see also Remark \\ref{rem:cont}).\nFirst note that by a result of Kechris, Solecki and Todor\\v{c}evi\\'c \\cite{KST} the \\emph{maximal independent set problem} (with any parameter $r\\in \\mathbb{N}$) is in this class for any bounded degree graph.\\footnote{That is, it is possible to find a Borel maximal independent set, i.e., a maximal independent set which is, moreover, a Borel subset of the vertex set.}\nIn particular, this yields that $\\mathsf{BOREL}$ contains the class $\\mathsf{LOCAL}(O(\\log^* n))$ by the characterization of \\cite{chang_kopelowitz_pettie2019exp_separation}, see \\cite{Bernshteyn2021LLL}. Moreover, as mentioned before, $\\mathsf{BOREL}$ is closed under countably many iterations of the operations of finding maximal independent set (for some parameter that might grow) and of applying a constant local rule that takes into account what has been constructed previously.\\footnote{It is in fact an open problem, whether this captures fully the class $\\mathsf{BOREL}$. However, note that an affirmative answer to this question would yield that problems can be solved in an ``effective\" manner in the Borel context, which is known not to be the case in unbounded degree graphs \\cite{todorvcevic2021complexity}.}\n\nTo get a better grasp of what this means, consider for example the proper vertex $2$-coloring problem on half-lines. It is clear that no local algorithm can solve this problem.\nHowever, as it is possible to determine the starting vertex after countably many iterations of the maximal independent set operation, we conclude that this problem is in the class $\\mathsf{BOREL}$.\nThe idea that $\\mathsf{BOREL}$ can compute some unbounded, global, information will be implicitly used in all the construction in \\cref{sec:separating_examples} that separate $\\mathsf{BOREL}$ from local classes.\n\nThe intuition behind the class $\\mathsf{MEASURE}$ is that it relates in the same way to the class $\\mathsf{BOREL}$, as randomized local algorithms relate to deterministic ones.\nIn particular, the operations that are allowed in the class $\\mathsf{MEASURE}$ are the same as in the class $\\mathsf{BOREL}$ but the solution of a given LCL can be incorrect on a measure zero set.\n\nThe class $\\mathsf{BAIRE}$ can be considered as a topological equivalent of the measure theoretic class $\\mathsf{MEASURE}$, that is, a solution can be incorrect on a topologically negligible set.\nThe main difference between the classes $\\mathsf{MEASURE}$ and $\\mathsf{BAIRE}$ is that in the later there is a hierarchical decomposition that is called \\emph{toast}.\n(Note that this phenomenon is present in the case of $\\mathsf{MEASURE}$ exactly on so-called amenable graphs. It is also tightly connected with the notion of hyperfiniteness \\cite{conleymillerbound,gao2015forcing}.)\nThe independence of colorings on a tree together with this structure allows for a combinatorial characterization of the class $\\mathsf{BAIRE}$, which was proven by Bernshteyn \\cite{Bernshteyn_work_in_progress}, see also \\cref{sec:baire}.\n\n\n\n\nNext we formulate the precise definitions.\nWe refer the reader to \\cite{pikhurko2021descriptive_comb_survey,kechris_marks2016descriptive_comb_survey,Bernshteyn2021LLL,kechrisclassical}, or to \\cite[Introduction,~Section~4.1]{grebik_rozhon2021LCL_on_paths} and \\cite[Section~7.1,~7.2]{grebik_rozhon2021toasts_and_tails} for intuitive examples and standard facts of descriptive set theory.\nIn particular, we do not define here the notion standard Borel\/probability space, a Polish topology, a Borel probability measure, Baire property etc.\n\nLet $\\mathcal{G}$ be a Borel graph of bounded maximum degree on a standard Borel space $X$.\nIn this paper we consider exclusively acyclic $\\Delta$-regular Borel graphs and we refer to them as \\emph{$\\Delta$-regular Borel forests}.\nIt is easy to see that the set of half-edges (see Definition \\ref{def:half-edge}) is naturally a standard Borel space, we denote this set by $H(\\mathcal{G})$.\nThus, it makes sense to consider Borel labelings of $H(\\mathcal{G})$.\nMoreover, if $\\mathcal{G}$ is a $\\Delta$-regular Borel forest and $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ is an LCL, we can also decide whether a coloring $f:\\mathcal{H}(\\mathcal{G})\\to \\Sigma$ is a solution to $\\Pi$ as in Definition \\ref{def:LCL_trees}. \nSimilarly, we say that the coloring $f$ solves $\\Pi$, e.g., on a $\\mu$-conull set for some Borel probability measure $\\mu$ on $X$ if there is a Borel set $C\\subseteq X$ such that $\\mu(C)=1$, the vertex constraints are satisfied around every $x\\in C$ and the edge constraints are satisfied for every $x,y\\in C$ that form an edge in $\\mathcal{G}$.\n\n\\begin{definition}[Descriptive classes]\n\\label{def:descriptive}\nLet $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ be an LCL.\nWe say that $\\Pi$ is in the class $\\mathsf{BOREL}$ if for every acyclic $\\Delta$-regular Borel graph $\\mathcal{G}$ on a standard Borel space $X$, there is a Borel function $f:H(\\mathcal{G})\\to \\Sigma$ that is a $\\Pi$-coloring of $\\mathcal{G}$.\n\nWe say that $\\Pi$ is in the class $\\mathsf{BAIRE}$ if for every acyclic $\\Delta$-regular Borel graph $\\mathcal{G}$ on a standard Borel space $X$ and every compatible Polish topology $\\tau$ on $X$, there is a Borel function $f:H(\\mathcal{G})\\to \\Sigma$ that is a $\\Pi$-coloring of $\\mathcal{G}$ on a $\\tau$-comeager set.\n\nWe say that $\\Pi$ is in the class $\\mathsf{MEASURE}$ if for every acyclic $\\Delta$-regular Borel graph $\\mathcal{G}$ on a standard Borel space $X$ and every Borel probability measure $\\mu$ on $X$, there is a Borel function $f:H(\\mathcal{G})\\to \\Sigma$ that is a $\\Pi$-coloring of $\\mathcal{G}$ on a $\\mu$-conull set.\n\\end{definition}\n\nThe following claim follows directly from the definition.\n\n\\begin{claim}\\label{cl:basicDC}\nWe have $\\mathsf{BOREL}\\subseteq \\mathsf{MEASURE}, \\mathsf{BAIRE}$.\n\\end{claim}\n\nRecently Bernshteyn \\cite{Bernshteyn2021LLL,Bernshteyn2021local=cont} proved several results that connect distributed computing with descriptive combinatorics.\nUsing the language of complexity classes we can formulate some of the result as inclusions in \\cref{fig:big_picture_trees}.\n\n\\begin{theorem}[\\cite{Bernshteyn2021LLL}]\nWe have\n\\begin{itemize}\n \\item $\\mathsf{LOCAL}(O(\\log^*(n)))\\subseteq \\mathsf{BOREL}$,\n \\item $\\mathsf{RLOCAL}(o(\\log(n))) \\subseteq \\mathsf{MEASURE}, \\mathsf{BAIRE}$.\n\\end{itemize}\nIn fact, these inclusions hold for any reasonable class of bounded degree graphs.\n\\end{theorem}\n\n\\begin{remark}\n\\label{rem:cont}\nThe ``truly'' local class in descriptive combinatorics is the class $\\mathsf{CONTINUOUS}$.\nEven though we do not define this class here\\footnote{To define the class $\\mathsf{CONTINUOUS}$, rather than asking for a continuous solution on all possible $\\Delta$-regular Borel graphs, one has to restrict to a smaller family of graphs, otherwise the class trivializes. To define precisely this family is somewhat inconvenient, and not necessary for our arguments.}, and we refer the reader to \\cite{Bernshteyn2021local=cont,GJKS,grebik_rozhon2021toasts_and_tails} for the definition and discussions in various cases, we mention that the inclusion \\ \\[\\mathsf{LOCAL}(O(\\log^* n))\\subseteq \\mathsf{CONTINUOUS} \\tag{*}\\] holds in most reasonable classes of bounded degree graphs, see \\cite{Bernshteyn2021LLL}.\nThis also applies to our situation\nRecently, it was shown by Bernshteyn \\cite{Bernshteyn2021local=cont}, and independently by Seward \\cite{Seward_personal}, that $(*)$ can be reversed for Cayley graphs of finitely generated groups.\nThis includes, e.g., natural $\\Delta$-regular trees with proper edge $\\Delta$-coloring, as this is a Cayely graph of free product of $\\Delta$-many $\\mathbb{Z}_2$ induced by the standard generating set.\nIt is, however, not clear whether it $(*)$ can be reversed in our situation, i.e., $\\Delta$-regular trees without any additional labels.\n\n\n\n\\end{remark}\n\n\n\n\n\n\n\\subsection{Random processes}\n\nWe start with an intutitve description of fiid processes. Let $T_\\Delta$ be the infinite $\\Delta$-regular tree.\nInformally, \\emph{factors of iid processes (fiid)} on $T_\\Delta$ are infinite analogues of local randomized algorithms in the following way.\nLet $\\Pi$ be an LCL and $u\\in T_\\Delta$.\nIn order to solve $\\Pi$, we are allowed to explore the whole graph, and the random strings that are assigned to vertices, and then output a labeling of half-edges around $u$.\nIf such an assignment is a measurable function and produces a $\\Pi$-coloring almost surely, then we say that $\\Pi$ is in the class $\\mathsf{fiid}$.\nMoreover, if every vertex needs to explore only finite neighborhood to output a solution, then we say that $\\Pi$ is in the class $\\mathsf{ffiid}$.\nSuch processes are called \\emph{finitary} fiid (\\emph{ffiid}). There is also an intimate connection between ffiid and uniform algorithms.\nThis is explained in \\cite[Section 2.2]{grebik_rozhon2021toasts_and_tails}.\nInformally, an ffiid process that almost surely solves $\\Pi$ is, in the language of distributed computing, an uniform local algorithm that solves $\\Pi$.\nThis allows us to talk about uniform local complexity of an ffiid.\nIn the rest of the paper we interchange both notions freely with slight preference for the distributed computing language.\n\nNow we define formally these classes, using the language of probability.\nWe denote by $\\operatorname{Aut}(T_\\Delta)$ the automorphism group of $T_\\Delta$.\nAn \\emph{iid process} on $T_\\Delta$ is a collection of iid random variables $Y=\\{Y_v\\}_{v\\in V(T_\\Delta)}$ indexed by vertices, or edges, half-edges etc, of $T_\\Delta$ such that their joint distribution is invariant under $\\operatorname{Aut}(T_\\Delta)$.\nWe say that $X$ is a \\emph{factor of iid process (fiid)} if $X=F(Y)$, where $F$ is a measurable $\\operatorname{Aut}(T_\\Delta)$-invariant map and $Y$ is an iid process on $T_\\Delta$.\\footnote{We always assume $Y=(2^\\mathbb{N})^{T_\\Delta}$ endowed with the product Lebesgue measure.}\nMoreover, we say that $X$ is a \\emph{finitary factor if iid process (ffiid)} if $F$ depends with probability $1$ on a finite (but random) radius around each vertex.\nWe denote as $R_F$ the random variable that assigns minimal such radius to a given vertex, and call it the \\emph{coding radius} of $F$.\nWe denote the space of all $\\Pi$-colorings of $T_\\Delta$ as $X_\\Pi$.\nThis is a subspace of $\\Sigma^{H(T_\\Delta)}$ that is invariant under $\\operatorname{Aut}(T_\\Delta)$.\n\n\\begin{definition}\nWe say that an LCL $\\Pi$ is in the class $\\mathsf{fiid}$ ($\\mathsf{ffiid}$) if there is an fiid (ffiid) process $X$ that almost surely produces elements of $X_\\Pi$.\n\nEquivalently, we can define $\\mathsf{ffiid} = \\bigcup_f \\mathsf{OLOCAL}(f(\\varepsilon))$ where $f$ ranges over all functions. That is, $\\mathsf{ffiid}$ is the class of problems solvable by any uniform distributed algorithm. \n\n\\end{definition}\n\n\nIt is obvious that $\\mathsf{ffiid}\\subseteq \\mathsf{fiid}$.\nThe natural connection between descriptive combinatorics and random processes is formulated by the inclusion $\\mathsf{MEASURE}\\subseteq \\mathsf{fiid}$.\nWhile this inclusion is trivially satisfied, e.g., in the case of $\\Delta$-regular trees with proper edge $\\Delta$-coloring, in our situation we need a routine argument that we include for the sake of completeness in \\cref{app:ktulu}.\n\n\\begin{restatable}{lemma}{FIIDinMEASURE}\\label{thm:FIIDinMEASURE}\nLet $\\Pi$ be an LCL such that $\\Pi\\in \\mathsf{MEASURE}$.\nThen $\\Pi\\in \\mathsf{fiid}$.\n\\end{restatable}\n\n\n\n\n\\subsection{Specific Local Problems}\n\nHere we list some well-known local problems that were studied in all three considered areas. We describe the results known about these problems with respect to classes from \\cref{fig:big_picture_trees}. \n\n\\paragraph{Edge Grabbing}\nWe start by recalling the well-known problem of edge grabbing (a close variant of the problem is known as sinkless orientation\\cite{brandt_etal2016LLL}). \nIn this problem, every vertex should mark one of its half-edges (that is, grab an edge) in such a way that no edge can be marked by two vertices. \nIt is known that $\\Pi_{\\text{edgegrab}} \\not\\in \\mathsf{LOCAL}(O(\\log^* n))$ \\cite{brandt_etal2016LLL} but $\\Pi_{\\text{edgegrab}} \\in \\mathsf{RLOCAL}(O(\\log\\log n))$.\n\nSimilarly, $\\Pi_{\\text{edgegrab}} \\not \\in \\mathsf{BOREL}$ by \\cite{DetMarks}, but $\\Pi_{\\text{edgegrab}}\\in \\mathsf{MEASURE}$ by \\cite{BrooksMeas}: to see the former, just note that if a $\\Delta$-regular tree admits proper $\\Delta$-colorings of both edges and vertices, every vertex can grab an edge of the same color as the color of the vertex. Thus, $\\Pi_{\\text{edgegrab}}\\in \\mathsf{BOREL}$ would yield that $\\Delta$-regular Borel forests with Borel proper edge-colorings admit a Borel proper $\\Delta$-coloring, contradicting \\cite{DetMarks}.\n\nThis completes the complexity characterization of $\\Pi_{\\text{edgegrab}}$ as well as the proper vertex $\\Delta$-coloring with respect to classes in \\cref{fig:big_picture_trees}. \n\n\n\n\n\n\\paragraph{Perfect Matching}\nAnother notorious LCL problem, whose position in \\cref{fig:big_picture_trees} is, however, not completely understood, is the perfect matching problem $\\Pi_{\\text{pm}}$.\nRecall that the perfect matching problem $\\Pi_{\\text{pm}}$ asks for a matching that covers all vertices of the input tree.\\footnote{In our formalism this means that around each vertex exactly one half-edge is picked.}\nIt is known that $\\Pi_{\\text{pm}} \\in \\mathsf{fiid}$ \\cite{lyons2011perfect}, and it is easy to see that $\\Pi_{\\text{pm}} \\not\\in \\mathsf{RLOCAL}(O(\\log\\log(n)))$ (we will prove a stronger result in \\cref{pr:AddingLineLLL}).\nMarks proved \\cite{DetMarks} that it is not in $\\mathsf{BOREL}$, even when the underlying tree admits a Borel bipartition.\nIt is not clear if $\\Pi_{\\text{pm}}$ is in $\\mathsf{ffiid}$, nor whether it is in $\\mathsf{MEASURE}$.\n\n\n\n\n\n\\paragraph{Graph Homomorphism}\nWe end our discussion with LCLs that correspond to graph homomorphisms (see also the discussion in Section \\ref{sec:marks}).\nThese are also known as edge constraint LCLs.\nLet $G$ be a finite graph.\nThen we define $\\Pi_G$ to be the LCL that asks for a homomorphism from the input tree to $G$, that is, vertices are labeled with vertices of $G$ in such a way that edge relations are preserved.\nThere are not many positive results except for the class $\\mathsf{BAIRE}$.\nIt follows from the result of Bernshteyn \\cite{Bernshteyn_work_in_progress} (see \\cref{sec:baire}) that $\\Pi_G\\in \\mathsf{BAIRE}$ if and only if $G$ is not bipartite.\nThe main negative results can be summarized as follows.\nAn easy consequence of the result of Marks \\cite{DetMarks} is that if $\\chi(G)\\le \\Delta$, then $\\Pi_G\\not \\in \\mathsf{BOREL}$.\nIn this paper, we describe examples of graphs of chromatic number up to $2\\Delta-2$ such that the corresponding homomorphism problem is not in $\\mathsf{BOREL}$, see \\cref{sec:marks}.\nThe theory of \\emph{entropy inequalities} see \\cite{backhausz} implies that if $G$ is a cycle on more than $9$ vertices, then $\\Pi_G\\not\\in \\mathsf{fiid}$.\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Generalization of Marks' technique}\n\\label{sec:marks}\n\nIn this section, we first develop a new way of proving lower bounds in the $\\mathsf{LOCAL}$ model based on a generalization of a technique of Marks \\cite{DetMarks}. Then, we use ideas arising in the finitary setting---connected to the adaptation of Marks' technique---to obtain new results back in the Borel context.\nFor an introduction to Marks' technique and a high-level discussion about the challenges in adapting the technique to the standard distributed setting as well as our solution via the new notion of an ID graph, we refer the reader to \\cref{sec:intromarks}.\n\n\\begin{comment}\nIn \\cite{DetMarks}, Marks showed that there exist $\\Delta$-regular Borel forests that do not admit a Borel proper $\\Delta$-coloring. As it was indicated above, this is a striking difference between $\\mathsf{MEASURE}$ or $\\mathsf{BAIRE}$ and $\\mathsf{BOREL}$, as such colorings exist in the former classes, see \\cite{BrooksMeas}. Since prior to Marks' breakthrough (essentially) any lower bound on Borel chromatic numbers utilized measure theoretic or Baire category arguments, his approach was (and had to be) based on a new technique: \\emph{determinacy of infinite games}. \n\nFix a payoff set $A \\subseteq \\mathbb{R}$. Consider the infinite game in which two players alternatingly determine digits of a real number. The first player wins iff the resulting real is in $A$. A rather counter-intuitive fact of set theory is that there exists a set $A$ for which none of the players has a winning strategy, i.e., the game is not determined. However, a deep theorem of Martin \\cite{martin} says that if $A$ is Borel, then one of the players does have such a strategy. \nMarks has constructed $\\Delta$-regular Borel forests and games with Borel payoff sets, for which the existence of a winning strategy for either of the players contradicts the existence of a Borel proper $\\Delta$-coloring. \n\nIn what follows, we generalize and adapt his games to the finitary setting, thus obtaining new lower bounds for LOCAL algorithms; then we use ideas inspired by the LOCAL technique of ID graphs to prove a new theorem in the Borel context. \n\nNote that Marks considered Borel versions of certain locally checkable problems (e.g., coloring, edge grabbing, perfect matching). As discussed above, for these problems, lower bounds were independently obtained via the round elimination technique in the LOCAL context, see e.g., \\cite{brandt19automatic}\\todo{cite more stuff??-Z}. Remarkably, the core ideas behind Brandt's and Marks' argument exhibit a surprising degree of similarity. \n\n\\todo{add the story about going from borel to local, and then back from local to borel}\n\nHowever, our extension of Marks' technique yields lower bounds for fairly natural problems that were resistant to round elimination.\nMoreover, as with our approach the lower bound is obtained by checking whether the considered LCL satisfies a certain combinatorial property called \\emph{playability condition} (see Definition~\\ref{def:Playability}), we can use it to obtain lower bounds for problems that have complicated descriptions---something that the round elimination technique in its current form is ill-suited for. \\todo{Maybe this sentence should rather be somewhere in the intro?}\n\n\n\n\\end{comment}\n\n\n\n\n\n\n\nThe setting we consider in this section is $\\Delta$-regular trees that come with a proper $\\Delta$-edge coloring with colors from $[\\Delta]$.\nAll lower bounds discussed already hold under these restrictions (and therefore also in any more general setting).\n\nRecall that an LCL $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ is given by specifying a list of allowed vertex configurations and a list of allowed edge configurations (see Definition \\ref{def:LCL_trees}).\nTo make our lower bounds as widely applicable as possible, we replace the latter list by a separate list for each of the $\\Delta$ edge colors; in other words, we consider LCLs where the correctness of the output is allowed to depend on the input that is provided by the edge colors.\nHence, formally, in this section, an LCL is more generally defined: it is a tuple $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ where $\\Sigma$ and $\\mathcal{V}$ are as before, while $\\mathcal{E} = (\\mathcal{E}_{\\alpha})_{\\alpha \\in [\\Delta]}$ is now a $\\Delta$-tuple of sets $\\mathcal{E}_{\\alpha}$ consisting of cardinality-$2$ multisets.\nSimilarly as before, a half-edge labeling (see Definition \\ref{def:descriptive} for the Borel context) with labels from $\\Sigma$ is a correct solution for $\\Pi$ if, for each vertex $v$, the multiset of labels assigned to the half-edges incident to $v$ is contained in $\\mathcal{V}$, and, for each edge $e$, the multiset of labels assigned to the half-edges belonging to $e$ is contained in $\\mathcal{E}_{\\alpha}$, where $\\alpha$ is the color of the edge $e$.\n\nThe idea of our approach is to identify a condition that an LCL $\\Pi$ necessarily has to satisfy if it is solvable in $O(\\log^* n)$ rounds in the $\\mathsf{LOCAL}$ model.\nShowing that $\\Pi$ does not satisfy this condition then immediately implies an $\\Omega(\\log n)$ deterministic and $\\Omega(\\log \\log n)$ randomized lower bound, by Theorem~\\ref{thm:basicLOCAL}.\n\nIn order to define our condition we need to introduce the notion of a configuration graph: a \\emph{configuration graph} is a $\\Delta$-tuple $\\mathbb{P}=(\\mathbb{P}_\\alpha)_{\\alpha\\in \\Delta}$ of graphs, where the vertex set of each of the graphs $\\mathbb{P}_\\alpha$ is the set of subsets of $\\Sigma$, and there is an edge in $\\mathbb{P}_\\alpha$ connecting two vertices $S,T$ if and only if there are $\\mathtt{a}\\in S$ and $\\mathtt{b}\\in T$ such that $\\{\\mathtt{a},\\mathtt{b}\\}\\in {\\mathcal{E}}_\\alpha$, note that loops are allowed.\n(Naturally, we will consider any two vertices of different $\\mathbb{P}_\\alpha$ to be distinct, even if they correspond to the same subset of $\\Sigma$.)\n\nNow we are set to define the aforementioned condition. The intuition behind the playability condition is the following: assume that there exists a local algorithm $\\mathcal{A}$ that solves $\\Pi$ using the $t$ neighbourhood of a given vertex. We are going to define a family of two player games. The game will depend on some $S \\subseteq \\Sigma$. $\\operatorname{Alice}$ and $\\operatorname{Bob}$ will assign labels (or IDs) to the vertices in the $t$-neighbourhood (in some way specified later on, depending on $\\alpha \\in [\\Delta]$). When the assignment is complete, we evaluate $\\mathcal{A}$ on the root of the obtained labelled graph, this way obtaining an element $s$ of $\\Sigma$, and decide who is the winner based on $s \n\\in S$ or not. Naturally, it depends on $S$ and $\\alpha$, which player has a winning strategy. This gives rise to a two coloring of vertices of $\\mathbb{P}_\\alpha$ by colors $\\operatorname{Alice}$ and $\\operatorname{Bob}$. The failure of the playability condition will guarantee that using a strategy stealing argument one can derive a contradiction.\n\n\n\n\n\\begin{definition}[Playability condition]\\label{def:Playability}\nWe say that an LCL $\\Pi$ is \\emph{playable} if for every $\\alpha\\in [\\Delta]$ there is a coloring $\\Lambda_{\\alpha}$ of the vertices of $\\mathbb{P}_{\\alpha}$ with two colors $\\{\\operatorname{Alice},\\operatorname{Bob}\\}$ such that the following conditions are satisfied:\n\\begin{enumerate}\n \\item [(A)] For any tuple $(S_{\\alpha})_{\\alpha \\in [\\Delta]} \\in V(\\mathbb{P}_1) \\times \\dots \\times V(\\mathbb{P}_{\\Delta})$ satisfying $\\Lambda_\\alpha(S_\\alpha)=\\operatorname{Alice}$ for each $\\alpha \\in [\\Delta]$, there exist $\\mathtt{a}_\\alpha\\not\\in S_\\alpha$ such that $\\{\\mathtt{a}_\\alpha\\}_{\\alpha\\in \\Delta}\\in {\\mathcal{V}}$, and\n \\item [(B)] for any $\\alpha \\in [\\Delta]$, and any tuple $(S,T) \\in V(\\mathbb{P}_{\\alpha}) \\times V(\\mathbb{P}_{\\alpha})$ satisfying $\\Lambda_\\alpha(S)=\\Lambda_\\alpha(T)=\\operatorname{Bob}$, we have that $(S,T)$ is an edge of $\\mathbb{P}_{\\alpha}$.\n\\end{enumerate}\n\n\\end{definition}\n\nOur aim in this section is to show the following general results.\n\n\\begin{theorem}\\label{th:MainLOCAL}\nLet $\\Pi$ be an LCL that is not playable.\nThen $\\Pi$ is not in the class $\\mathsf{LOCAL}(O(\\log^* n))$.\n\\end{theorem}\n\nUsing ideas from the proof of this result, we can formulate an analogous theorem in the Borel context.\nLet us mention that while \\cref{th:MainLOCAL} is a consequence of \\cref{thm:MainBorel} by \\cite{Bernshteyn2021LLL}, we prefer to state the theorems separately.\nThis is because the proof of the Borel version of the theorem uses heavily complex results such as the Borel determinacy theorem, theory of local-global limits and some set-theoretical considerations about measurable sets.\nThis is in a stark contrast with the proof of \\cref{th:MainLOCAL} that uses `merely' the existence of large girth graphs and determinacy of finite games.\nHowever, the ideas surrounding these concepts are the same in both cases.\n\n\\begin{theorem}\\label{thm:MainBorel}\nLet $\\Pi$ be an LCL that is not playable\nThen $\\Pi$ is not in the class $\\mathsf{BOREL}$.\n\\end{theorem}\n\nThe main application of these abstract results is to find lower bounds for \\emph{(graph) homomorphism LCLs} aka \\emph{edge constraint LCLs}.\nWe already defined this class in \\cref{subsec:local_problems} but we recall the definition in the formalism of this section to avoid any confusion.\nLet $G$ be a finite graph.\nThen $\\Pi_G=(\\Sigma,\\mathcal{V},\\mathcal{E})$ is defined by letting\n\\begin{enumerate}\n\\item \\label{p:vertex} ${\\mathcal{V}}=\\{({\\mathtt{a}},\\dots,{\\mathtt{a}}):\\mathtt{a}\\in \\Sigma\\}$,\n\\item $\\Sigma=V(G)$,\n\\item $\\forall \\alpha \\in [\\Delta] \\ (\\mathcal{E}_\\alpha=E(G)).$\n\\end{enumerate}\nAn LCL for which \\eqref{p:vertex} holds is called a \\emph{vertex-LCL}. There is a one-to-one correspondence between vertex-LCLs for which $\\forall \n\\alpha,\\beta \\ (\\mathcal{E}_\\alpha=\\mathcal{E}_\\beta)$ and LCLs of the form $\\Pi_G$. (Indeed, to vertex-LCLs with this property one can associate a graph $G$ whose vertices are the labels in $\\Sigma$ and where two vertices $\\ell, \\ell'$ are connected by an edge if and only if $\\{ \\ell, \\ell' \\} \\in \\mathcal{E}$.)\nNote that if $\\Pi$ is a vertex-LCL, then condition (A) in Definition~\\ref{def:Playability} is equivalent to the statement that no tuple $\\{S_\\alpha\\}_{\\alpha\\in \\Delta}$ that satisfies the assumption of (A) can cover $\\Sigma$, i.e., that $\\Sigma \\nsubseteq S_1 \\cup \\dots \\cup S_{\\Delta}$ for such a tuple.\nMoreover, if $\\Pi_G$ is a homomorphism LCL, then $\\mathbb{P}_\\alpha=\\mathbb{P}_{\\beta}$ for every $\\alpha,\\beta\\in \\Delta$.\n\nThe simplest example of a graph homomorphism LCL is given by setting $G$ to be the clique on $k$ vertices; the associated LCL $\\Pi_G$ is simply the problem of finding a proper $k$-coloring of the vertices of the input graph.\nNow we can easily derive Marks' original result from our general theorem.\n\n\\begin{corollary}[Marks \\cite{DetMarks}]\\label{cor:detmarks}\nLet $G$ be a finite graph that has chromatic number at most $\\Delta$. Then $\\Pi_G$ is not playable. In particular, there is a $\\Delta$-regular Borel forest that have Borel chromatic number $\\Delta+1$. \n\\end{corollary}\n\\begin{proof}\nLet $A_1,\\dots, A_\\Delta$ be some independent sets that cover $G$, and $\\Lambda_1, \\dots, \\Lambda_{\\Delta}$ arbitrary colorings of the vertices of $\\mathbb{P}_1, \\dots, \\mathbb{P}_{\\Delta}$, respectively, with colors from $\\{\\operatorname{Alice},\\operatorname{Bob}\\}$.\nIt follows that $\\Lambda_\\alpha(A_\\alpha)=\\operatorname{Alice}$ for every $\\alpha\\in \\Delta$, since otherwise condition (B) in Definition~\\ref{def:Playability} is violated with $S=T=A_\\alpha$.\nBut then condition (A) does not hold.\n\\end{proof}\n\nAs our main application we describe graphs with chromatic number larger than $\\Delta$ such that $\\Pi_G$ is not playable.\nThis rules out the hope that the complexity of $\\Pi_G$ is connected with chromatic number being larger than $\\Delta$.\nIn Section \\ref{subsec:homlcls} we show the following.\nNote that the case $k=\\Delta$ is \\cref{cor:detmarks}.\n\n\\begin{theorem}\n\\label{th:homomorphism}\nLet $\\Delta > 2$ and $\\Delta < k \\leq 2\\Delta-2$. There exists a graph $G_k$ with $\\chi(G_k)= k$, such that $\\Pi_{G_k}$ is not playable and $\\Pi_{G_k} \\in \\mathsf{RLOCAL}(O(\\log\\log n ))$. In particular, $\\Pi_{G_k}\\not\\in \\mathsf{BOREL}$ and $\\Pi_{G_k}\\not\\in\\mathsf{LOCAL}(O(\\log^*(n)))$.\n\\end{theorem}\n\nInterestingly, recent results connected to counterexamples to Hedetniemi's conjecture yield the same statement asymptotically, as $\\Delta \\to \\infty$ (see Remark \\ref{r:hedet}).\n\n\\begin{remark}\nIt can be shown that for $\\Delta=3$ both the Chv\\' atal and Gr\\\" otsch graphs are suitable examples for $k=4$.\n\\end{remark}\n\n\\begin{remark}\nAs another application of his technique Marks showed in \\cite{DetMarks} that there is a Borel $\\Delta$-regular forest that does not admit Borel perfect matching.\nThis holds even when we assume that the forest is Borel bipartite, i.e., it has Borel chromatic number $2$.\nIn order to show this result Marks works with free joins of colored hyperedges, that is, Cayley graphs of free products of cyclic groups.\nOne should think of two types of triangles ($3$-hyperedges) that are joined in such a way that every vertex is contained in both types and there are no cycles.\nWe remark that the playability condition can be formulated in this setting.\nSimilarly, one can derive a version of \\cref{thm:MainBorel}.\nHowever, we do not have any application of this generalization.\n\\end{remark}\n\n\n\n\n\n\n\n\n\\subsection{Applications of playability to homomorphism LCLs}\\label{subsec:homlcls}\n\nIn this section we find the graph examples from \\cref{th:homomorphism}.\nFirst we introduce a condition $\\Delta$-(*) that is a weaker condition than having chromatic number at most $\\Delta$, but still implies that the homomorphism LCL is not playable. Then, we will show that--similarly to the way the complete graph on $\\Delta$-many vertices, $K_\\Delta$, is maximal among graphs of chromatic number $\\leq \\Delta$--there exists a maximal graph (under homorphisms) with property $\\Delta$-(*). Recall that we assume $\\Delta >2$.\n\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=.6\\textwidth]{fig\/h3.pdf}\n \\caption{The maximal graph with the property $\\Delta$-(*) for $\\Delta=3$.\n \\label{fig:h3}\n\\end{figure}\n\n\n\n\n\\begin{definition}[Property $\\Delta$-(*)]\nLet $\\Delta>2$ and $G=(V,E)$ be a finite graph.\nWe say that $G$ satisfies property $\\Delta$-(*) if there are sets $S_0,S_1\\subseteq V$ such that $G$ restricted to $V \\setminus S_i$ has chromatic number at most $(\\Delta-1)$ for $i\\in \\{0,1\\}$, and there is no edge between $S_0$ and $S_1$.\n\\end{definition}\n\nNote that $\\chi(G) \\leq \\Delta$ implies $\\Delta$-(*): indeed, if $A_1,\\dots,A_\\Delta$ are independent sets that cover $V(G)$, we can set $S_0=S_1=A_1$.\n\nOn the other hand, we claim that if $G$ satisfies $\\Delta$-(*) then $\\chi(G)\\le 2\\Delta-2$. In order to see this, take $S_0,S_1 \\subseteq V(G)$ witnessing $\\Delta$-(*). Then, as there is no edge between $S_0$ and $S_1$, so in particular, between $S_0 \\setminus S_1$ and $S_0 \\cap S_1$, it follows that the chromatic number of $G$'s restriction to $S_0$ is $\\leq \\Delta-1$. But then we can construct proper $\\Delta-1$-colorings of $S_0$ and $V(G) \\setminus S_0$, which shows our claim.\n\n\n\\begin{proposition}\\label{pr:CONdition*}\nLet $G$ be a graph satisfying $\\Delta$-(*).\nThen $\\Pi_G$ is not playable.\n\\end{proposition}\n\\begin{proof}\nFix $S_0,S_1$ as in the definition of $\\Delta$-(*) and assume for a contradiction that colorings $\\Lambda_1, \\dots, \\Lambda_\\Delta$ as described in Definition~\\ref{def:Playability} exist.\nBy Property $\\Delta$-(*), there exist independent sets $A_1, \\dots, A_{\\Delta-1}$ such that $S_0$ together with the $A_i$ covers $G$, i.e., such that $S_0 \\cup A_1 \\cup \\dots \\cup A_{\\Delta-1} = V(G)$.\nFor each $\\alpha \\in [\\Delta-1]$, we must have $\\Lambda_\\alpha(A_{\\alpha})=\\operatorname{Alice}$, since otherwise condition (B) in Definition~\\ref{def:Playability} is violated.\nConsequently $\\Lambda_{\\Delta}(S_0)=\\operatorname{Bob}$, otherwise condition (A) is violated.\nSimilarly $\\Lambda_\\Delta(S_1)=\\operatorname{Bob}$.\nThis shows that $\\Lambda_\\Delta$ does not satisfy condition (B) with $S = S_0$ and $T = S_1$.\n\\end{proof}\n\nNext we describe maximal examples of graphs that satisfy the condition $\\Delta$-(*). That is, we define a graph $H_\\Delta$ that satisfies $\\Delta$-(*), its chromatic number is $2\\Delta-2$ and every other graph that satisfies $\\Delta$-(*) admits a homomorphism in $H_\\Delta$.\nThe graph $H_{3}$ is depicted in \\cref{fig:h3}.\n\nRecall that the (categorical) product $G \\times H$ of graphs $G,H$ is the graph on $V(G) \\times V(H)$, such that $((g,h),(g',h')) \\in E(G \\times H)$ iff $(g,g') \\in E(G)$ and $(h,h') \\in E(H)$.\n \n Write $P$ for the product $K_{\\Delta-1}\\times K_{\\Delta-1}$. Let $V_0$ and $V_1$ be vertex disjoint copies of $K_{\\Delta-1}$.\nWe think of vertices in $V_i$ and $P$ as having labels from $[\\Delta-1]$ and $[\\Delta-1] \\times [\\Delta-1]$, respectively.\nThe graph $H_\\Delta$ is the disjoint union of $V_0$, $V_1$, $P$ and an extra vertex $\\dagger$ that is connected by an edge to every vertex in $P$, and additionally, if $v$ is a vertex in $V_0$ with label $i\\in [\\Delta-1]$, then we connect it by an edge with $(i',j)\\in P$ for every $i'\\not=i$ and $j\\in [\\Delta-1]$, and if $v$ is a vertex in $V_1$ with label $j\\in \\Delta-1$, then we connect it by an edge with $(i,j')\\in P$ for every $j'\\not=j$ and $i \\in [\\Delta-1]$.\n\n\n\\begin{proposition}\\label{pr:MinimalExample}\n\\begin{enumerate}\n \\item \\label{prop:hdelta}$H_\\Delta$ satisfies $\\Delta$-(*).\n \\item \\label{prop:hdeltachrom} $\\chi(H_\\Delta)=2\\Delta-2$.\n \\item \\label{prop:hdeltamax} A graph $G$ satisfies $\\Delta$-(*) if and only if it admits a homomorphism to $H_\\Delta$.\n\\end{enumerate}\n\n\\end{proposition}\n\\begin{proof}\n\\eqref{prop:hdelta} Set $S_0=V(V_0)\\cup \\{\\dagger\\}$ and $S_1=V(V_1)\\cup \\{\\dagger\\}$.\nBy the definition there are no edges between $S_0$ and $S_1$.\nConsider now, e.g., $V(H_\\Delta)\\setminus S_0$.\nLet $A_j$ consist of all elements in $P$ that have second coordinate equal to $j$ together with the vertex in $V_1$ that has the label $j$.\nBy the definition, the set $A_i$ is independent and $\\bigcup_{i\\in [\\Delta-1]}A_i$ covers $H_\\Delta\\setminus S_0$, and similarly for $S_1$.\n\n\\eqref{prop:hdeltachrom} By \\eqref{prop:hdelta} and the claim after the definition of $\\Delta$-(*), it is enough to show that $\\chi(H_\\Delta)\\geq 2\\Delta-2$. Towards a contradiction, assume that $c$ is a proper coloring of $H_\\Delta$ with $<2\\Delta-2$-many colors. Note the vertex $\\dagger$ guarantees that $|c(V(P))|\\leq 2\\Delta-4$, and also $\\Delta-1 \\leq |c(V(P))|$.\n\nFirst we claim that there are no indices $i,j\\in [\\Delta-1]$ (even with $i=j$) such that $c(i,r)\\not=c(i,s)$ and $c(r,j)\\not=c(s,j)$ for every $s\\not=r$: indeed, otherwise, by the definition of $P$ we had $c(i,r)\\not=c(s,j)$ for every $r,s$ unless $(i,r)=(s,j)$, which would the upper bound on the size of $c(V(P))$.\n\nTherefore, without loss of generality, we may assume that for every $i\\in [\\Delta-1]$ there is a color $\\alpha_i$ and two indices $j_i \\neq j'_i$ such that $c(i,j_i)=c(i,j'_i)=\\alpha_i$. It follows form the definition of $P$ and $j_i \\neq j'_i$ that\n $\\alpha_i\\not=\\alpha_{i'}$ whenever $i\\not= i'$. \n \n Moreover, note that any vertex in $V_1$ is connected to at least one of the vertices $(i,j_i)$ and $(i,j'_i)$, hence none of the colors $\\{\\alpha_i\\}_{i \\in [\\Delta-1]}$ can appear on $V_1$. Consequently, since $V_1$ is isomorphic to $K_{\\Delta-1}$ we need to use at least $\\Delta-1$ additional colors, a contradiction.\n\n\\eqref{prop:hdeltamax} First note that if $G$ admits a homomorphism into $H_\\Delta$, then the pullbacks of the sets witnessing $\\Delta$-(*) will witness that $G$ has $\\Delta$-(*).\n\nConversely, let $G$ be a graph that satisfies $\\Delta$-(*).\nFix the corresponding sets $S_0,S_1$ together with $(\\Delta-1)$-colorings $c_0,c_1$ of their complements.\nWe construct a homomorphism $\\Theta$ from $G$ to $H_\\Delta$. Let\n\\[\\Theta(v)=\n\\begin{cases} \\dagger &\\mbox{if } v \\in S_0 \\cap S_1, \\\\\nc_0(v) & \\mbox{if } v \\in S_1 \\setminus S_0, \\\\\nc_1(v) & \\mbox{if } v \\in S_0 \\setminus S_1,\n\\\\\n(c_0(v),c_1(v)) & \\mbox{if } v \\not \\in S_0 \\cup S_1.\n\\end{cases} \\]\n\n\nObserve that $S=S_0\\cap S_1$ is an independent set such that there is no edge between $S$ and $S_0 \\cup S_1$.\nUsing this observation, one easily checks case-by-case that $\\Theta$ is indeed a homomorphism.\n\n \n \n \n \n \n \n\\end{proof}\n\nNow, combining what we have so far, we can easily prove Theorem \\ref{th:homomorphism}.\n\n\n\\begin{proof}[Proof of Theorem \\ref{th:homomorphism}]\nIt follows that $\\Pi_{H_\\Delta}$ is not playable from \\cref{pr:MinimalExample}, \\cref{pr:CONdition*}. It is easy to see that if for a graph $G$ the LCL $\\Pi_G$ is not playable then $\\Pi_{G'}$ is not playable for every subgraph $G'$ of $G$. Since erasing a vertex decreases the chromatic number with at most one, for each $k\\leq 2\\Delta-2$ there is a subgraph $G_k$ of $H_\\Delta$ with $\\chi(G_k)=k$, such that $\\Pi_{G_k}$ is not playable.\n\nIt follows from Theorems \\ref{th:MainLOCAL} and \\ref{thm:MainBorel} that there is a $\\Delta$-regular Borel forest that admits no Borel homomorphism to any graph of $G_k$ and that $\\Pi_{G_k} \\not \\in \\mathsf{LOCAL}(O(\\log^*(n)))$.\n\nFinally, note that if $k \\geq \\Delta$ then $G_k$ can be chosen so that it contains $K_\\Delta$, yielding $\\Pi_{G_k} \\in \\mathsf{RLOCAL}(O(\\log\\log(n))$.\n\\end{proof}\n\n\\begin{remark}\n\\label{r:hedet}Recall that Hedetniemi's conjecture is the statement that if $G,H$ are finite graphs then $\\chi(G \\times H)=\\min\\{\\chi(G),\\chi(H)\\}$. This conjecture has been recently disproven by Shitov \\cite{shitov}, and strong counterexamples have been constructed later (see, \\cite{tardif2019note,zhu2021note}). We claim that these imply for $\\varepsilon>0$ the existence of finite graphs $H$ with $\\chi(H) \\geq (2-\\varepsilon)\\Delta$ to which $\\Delta$-regular Borel forests cannot have a homomorphism in BOREL, for every large enough $\\Delta$. Indeed, if a $\\Delta$-regular Borel forest admitted a Borel homomorphism to each finite graph of chromatic number at least $(2-\\varepsilon)\\Delta$, it would have such a homomorphism to their product as well. Thus, we would obtain that the chromatic number of the product of any graphs of chromatic number $(2-\\varepsilon)\\Delta$ is at least $\\Delta+1$. This contradicts Zhu's result \\cite{zhu2021note}, which states that the chromatic number of the product of graphs with chromatic number $n$ can drop to $\\approx \\frac{n}{2}$.\n\n\\end{remark}\n\n\\begin{remark}\nA natural attempt to construct graphs with large girth and not playable homomorphisms problem would be to consider random $d$-regular graphs of size $n$ for a large enough $n$. However, it is not hard to see that setting $\\Lambda(A)=\\operatorname{Alice}$ if and only if $|A|<\\frac{n}{d}$ shows that this approach cannot work.\n\\end{remark}\n\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{Proof of \\cref{th:MainLOCAL}}\\label{sec:marksproof}\n\nIn this section we prove Theorem \\ref{th:MainLOCAL}, by applying Marks' game technique in the LOCAL setting. In order to define our games, we will need certain auxiliary graphs, the so-called \\emph{ID graphs}. \nThe purpose of these graphs is to define a ``playground'' for the games that we consider.\nNamely, vertices in the game are labeled by vertices from the ID graph in such a way that the at the end we obtain a homomorphism from our underlying graph to the ID graph. \n\n\n\n\n\\begin{definition}\nA pair $\\idgraphcolored{n}{t}{r}=(\\idgraph{n}{t}{r},c)$ is called an \\emph{ID graph}, if\n\\begin{enumerate}\n \\item \\label{p:idgirth} $\\idgraph{n}{t}{r}$ is graph with girth at least $2t+2$,\n \\item \\label{p:idsize} $|V(\\idgraph{n}{t}{r})|\\le n$,\n \\item \\label{p:idlabels} $c$ is a $\\Delta$-edge-coloring of $\\idgraph{n}{t}{r}$, such that every vertex is adjacent to at least one edge of each color,\n \\item \\label{p:idratio} for each $\\alpha \\in [\\Delta]$ the ratio of a maximal independent set of $\\idgraph{n}{t}{r}^{\\alpha}=(V(\\idgraph{n}{t}{r}), E(\\idgraph{n}{t}{r}) \\cap c^{-1}(\\alpha))$ is at most $r$ (i.e., $\\idgraph{n}{t}{r}^\\alpha$ is the graph formed by $\\alpha$-colored edges).\n\\end{enumerate}\n\\end{definition}\n\nBefore we define the game we show that ID graphs exist.\n\n\\begin{proposition}\\label{pr:IDgraph}\nLet $t_n\\in o(\\log(n))$, $r>0$, $\\Delta\\geq 2$.\nThen there is an ID graph $\\idgraphcolored{n}{t_n}{r}$ for every $n\\in \\mathbb{N}$ sufficiently large.\n\\end{proposition}\n\\begin{proof}\nWe use the \\emph{configuration model} for regular random graphs, see \\cite{wormald1999models}.\nThis model is defined as follows.\nLet $n$ be even.\nThen a $d$-regular random sample on $n$-vertices is just a union of $d$-many independent uniform random perfect matchings.\nNote that in this model we allow parallel edges.\n\nIt was proved by Bollob\\'as \\cite{bollobas} that the independence ratio of a random $d$-regular graph is at most $\\frac{2\\log(d)}{d}$ a.a.s.\nMoreover, this quantity is concentrated by an easy application of McDiarmid's result \\cite{mcdiarmid2002concentration}, i.e., \n\\begin{equation}\n \\P\\left(\\left|X-\\textrm{\\textbf{E}}(X)\\right|\\ge \\sqrt{n}\\right)<2\\exp\\left(-\\frac{n}{d}\\right),\n\\end{equation}\nwhere $X$ is the random variable that counts the size of a maximal independent set.\nTherefore for fixed $n$ large enough we have that the independence ratio of a random sample is at most $\\frac{3\\log(d)}{d}$ with probability at least $\\left(1-2\\exp\\left(-\\frac{n}{d}\\right)\\right)$.\n\n\nPick a $d$ large enough such that $\\frac{3\\log(d)}{d} & 1-\\left(1-2\\exp\\left(-\\frac{n}{d}\\right)\\right)^\\Delta\n\\end{split}\n\\end{equation}\nas $n\\to \\infty$.\nThis shows that there exists such a graph $\\idgraphcolored{t_n}{n}{r}$ with non-zero probability.\n\\end{proof}\n\n\n\n\nNext, we define the games. As mentioned before, the games are going to depend on the following parameters: an algorithm $\\mathcal{A}_n$ of local complexity $t\\in o(\\log(n))$, an ID graph $\\idgraphcolored{n}{t}{r}$, $\\alpha\\in \\Delta$, $\\sigma\\in V(\\idgraph{n}{t}{r})$ and $S \\subseteq \\Sigma$.\n(We will view $\\mathcal{A}_n$, $\\idgraphcolored{n}{t}{r}$ as fixed, and the rest of the parameters as variables).\n\nThe game\n$$\\mathbb{G}(\\mathcal{A}_n,n,t,\\idgraphcolored{n}{t}{r})[\\alpha,\\sigma,S]$$\nis defined as follows: two players, $\\operatorname{Alice}$ and $\\operatorname{Bob}$ assign labels to the vertices of a rooted $\\Delta$-regular tree of diameter $t$. The labels are vertices of $\\idgraph{n}{t}{r}$ and the root is labeled $\\sigma$. In the $k$-th round, where $00$. $\\idgraphingcolored{r}=(\\idgraphing{r},c)$ is an \\emph{ID graphing}, if\n\n\\begin{enumerate}\n \\item $\\idgraphing{r}$ is an acyclic locally finite Borel graphing on a standard probability measure space $(X,\\mu)$,\n \\item \\label{p:idgraphinglabels} $c$ is a Borel $\\Delta$-edge-coloring of $\\idgraphing{r}$, such that every vertex is adjacent to at least one edge of each color,\n \\item \\label{p:idgraphingratio} for each $\\alpha \\in [\\Delta]$ the $\\mu$-measure of a maximal independent set of $\\idgraphing{r}^{\\alpha}=(V(\\idgraphing{r}), E(\\idgraphing{r}) \\cap c^{-1}(\\alpha))$ is at most $r$.\n\\end{enumerate}\n\\end{definition}\n\n\\begin{proposition}\n\\label{p:idgraphing}\nFor each $r>0$ there exists an ID graphing $\\idgraphingcolored{r}$.\n\\end{proposition}\n\\begin{proof}\nLet $\\idgraphing{r}$ be a \\emph{local-global limit} of the random graphs constructed in \\cref{pr:IDgraph} (see, e.g., \\cite{hatamilovaszszegedy} for the basic results about local-global convergence).\nIt is not hard to check that this limit satisfies the required properties.\n\\end{proof}\n\n\nNow we are ready to prove the theorem. The proof will closely follow the argument given in the proof of the LOCAL version, i.e., Theorem \\ref{th:MainLOCAL}, but can be understood without reading the latter.\n\n\\begin{proof}[Proof of \\cref{thm:MainBorel}]\nLet $\\Pi$ be an LCL that is not playable, and $N\\in \\mathbb{N}$ be the number of all possible colorings of vertices of $(\\mathbb{P}_{\\alpha})_{\\alpha\\in \\Delta}$ with two colors, that is, $N=2^{\\sum_\\alpha |V(\\mathbb{P}_\\alpha)|}$.\nSet $r:=\\frac{1}{N+1}$.\n\nWe define a Borel acyclic $\\Delta$-regular graph $\\mathcal{G}$, with edges properly colored by $\\Delta$, that does not admit a Borel solution of $\\Pi$.\nVertices of $\\mathcal{G}$ are pairs $(x,A)$, where $x\\in X$ is a vertex of $\\idgraphing{r}$ and $A$ is a countable subgraph of $\\idgraphing{r}$ that is a $\\Delta$-regular tree that contains $x$ and the edge coloring of $\\idgraphingcolored{r}$ induces a proper edge coloring of $A$.\nWe say that $(x,A)$ and $(y,B)$ are connected by an $\\alpha$-edge in $\\mathcal{G}$ if $A=B$, $x,y$ are adjacent in $A$ and the edge that connects them has color $\\alpha\\in \\Delta$.\n\nSuppose for a contradiction that $\\mathcal{A}$ is a Borel function that solves $\\Pi$ on $\\mathcal{G}$.\n\nNext, we define a family of games parametrized by $\\alpha \\in [\\Delta]$, $x \\in V(\\idgraphing{r})$ and $S \\subseteq \\Sigma$. For the reader familiar with Marks' construction, let us point out that for a fixed $x$, the games are analogues to the ones he defines, with the following differences: allowed moves are vertices of the ID graphing $\\idgraphing{r}$ and restricted by its edge relation, and the winning condition is defined by a set of labels, not just merely one label.\n\nSo, the game\n$$\\mathbb{G}(\\mathcal{A},\\idgraphingcolored{r})[\\alpha,x,S]$$ \nis defined as follows: $\\operatorname{Alice}$ and $\\operatorname{Bob}$ alternatingly label vertices of a $\\Delta$-regular rooted tree. The root is labelled by $x$, and the labels come from $V(\\idgraphing{r})$. In the $k$-th round, first $\\operatorname{Alice}$ labels vertices of distance $k$ from the root on the side of the $\\alpha$ edge.\nAfter that, $\\operatorname{Bob}$ labels all remaining vertices of distance $k$, etc (see \\cref{fig:game}). We also require the assignment of labels to give rise to an edge-color preserving homomorphism to $\\idgraphingcolored{r}$. \n\nIt follows from the acyclicity of $\\idgraphing{r}$ that a play of a game determines a $\\Delta$-regular rooted subtree of $\\idgraphing{r}$ to which the restriction of the edge-coloring is proper. That is, it determines a vertex $(x,A)$ of $\\mathcal{G}$. Let $\\operatorname{Alice}$ win iff the output of $\\mathcal{A}$ on the half-edge determined by $(x,A)$ and the color $\\alpha$ is \\textbf{not} in $S$.\n\nDefine the function $\\Lambda^x_{\\alpha}:\\mathbb{P}_\\alpha\\to \\{\\operatorname{Alice},\\operatorname{Bob}\\}$ assigning the player to some $S \\in V(\\mathbb{P}_\\alpha)$ who has a winning strategy in $\\mathbb{G}(\\mathcal{A},\\idgraphingcolored{r})[\\alpha,x,S]$.\nNote that, since $\\idgraphing{r}$ is locally finite, each player has only finitely many choices at each position. Thus, it follows from Borel Determinacy Theorem that $\\Lambda^x_{\\alpha}$ is well defined.\n\nNow we show the analogue of Proposition \\ref{pr:(A)satisfied}.\n\n\\begin{proposition}\\label{pr:Borel(A)satisfied}\n$(\\Lambda^{x}_{\\alpha})_{\\alpha\\in \\Delta}$ satisfies (A) in \\cref{def:Playability}.\n\\end{proposition}\n\\begin{proof}\nAssume that $(S_\\alpha)_{\\alpha\\in \\Delta}$ is such that $\\Lambda^x_\\alpha(S_\\alpha)=\\operatorname{Alice}$.\nThis means that $\\operatorname{Alice}$ has winning strategy in all the games corresponding to $S_\\alpha$. Letting these strategies play against each other in the obvious way, produces a vertex $(x,A)$ of $\\mathcal{G}$.\nSince $\\mathcal{A}$ solves $\\Pi$ it has to output labeling of half edges $(\\mathtt{a}_{\\alpha})_{\\alpha\\in \\Delta}\\in \\mathcal{N}$, where $\\mathtt{a}_\\alpha$ is a label on the half edge that start at $(x,A)$ and has color $\\alpha$.\nNote that we must have $\\mathtt{a}_\\alpha\\not\\in S_\\alpha$ by the definition of winning strategy for $\\operatorname{Alice}$.\nThis shows that (A) of Definition \\ref{def:Playability} holds.\n\\end{proof}\n\n Let $f$ be \n defined by $$x\\mapsto (\\Lambda^x_{\\alpha})_{\\alpha\\in\\Delta}.$$ Note that $f$ has a finite range. Using the fact that the allowed moves for each player can be determined in a Borel way, uniformly in $x$, it is not hard to see that for each element $s$ in the range, $f^{-1}(s)$ is in the algebra generated by sets that can be obtained by applying game quantifiers to Borel sets (for the definition see \\cite[Section 20.D]{kechrisclassical}). It has been shown by Solovay \\cite{Solovay} and independently by Fenstad-Norman \\cite{fenstadnorman} that such sets are provably $\\Delta^1_2$, and consequently, measurable. Therefore, $f$ is a measurable map.\n \n As the number of sequences of functions $(\\Lambda^x_{\\alpha})_{\\alpha\\in\\Delta}$ is $\\leq N$, by the choice of $r$, there exists a Borel set $Y$ with $\\mu(Y) > r$, such that $f$ is constant on $Y$.\n \nSince $\\Pi$ is not playable Proposition \\ref{pr:Borel(A)satisfied} gives that there is an $\\alpha\\in \\Delta$ and $S,T\\in \\mathbb{P}_\\alpha$ that violate condition (B).\nRecall that the measure of an independent Borel set of $\\idgraphing{r}^\\alpha$ is at most $r$.\nThat means that there are $x,y\\in Y$ such that $(x,y)$ is an $\\alpha$-edge in $\\idgraphing{r}$.\nWe let the winning strategies of $\\operatorname{Bob}$ in the games\n$$\\mathbb{G}(\\mathcal{A},\\idgraphingcolored{r})[\\alpha,x,S], \\ \\mathbb{G}(\\mathcal{A},\\idgraphingcolored{r})[\\alpha,y,T],$$\nplay against each other, where we start with an $\\alpha$ edge with endpoints labeled by $x,y$.\nThis produces a labeling of a $\\Delta$-regular tree $A$ with labels from $\\idgraphing{r}$ such that $(x,A)$ and $(y,A)$ span an $\\alpha$-edge in $\\idgraphing{r}$.\nApplying $\\mathcal{A}$ on the half-edge determined by $(x,A)$ and the color $\\alpha$, we obtain $\\mathtt{a}_0\\in S$.\nSimilarly, we produce $\\mathtt{a}_1\\in T$ for $(y,A)$.\nHowever, $(\\mathtt{a}_0,\\mathtt{a}_1)\\not\\in {\\bf E}_\\alpha$ by the definition of an edge in $\\mathbb{P}_{\\alpha}$.\nThis shows that $\\mathcal{A}$ does not produce a Borel solution to $\\Pi$.\n\\end{proof}\n\n\\begin{remark}\nOne can give an alternative proof of the existence of a Borel $\\Delta$-regular forest that does not admit a Borel homomorphism to $G_k$ that avoids using the graphing $\\idgraphing{r}$ and uses the original example described by Marks \\cite{DetMarks} instead.\nThe reason is that the example graphs $G_k$ contain $K_\\Delta$, as discussed in the proof of \\cref{th:homomorphism}.\nThis takes care of the ``non-free'' as in the argument of Marks.\nThen it is enough to use the pigeonhole principle on $\\mathbb{N}$ instead of the independence ratio reasoning.\n\\end{remark}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Separation Of Various Complexity Classes}\n\\label{sec:separating_examples}\n\nIn this section we provide examples of several problems that separate some classes from \\cref{fig:big_picture_trees}.\nThe examples show two things. First, we have $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon) \\not= \\mathsf{OLOCAL}(O(\\log\\log 1\/\\varepsilon))$. \nThis shows that there are problems such that their worst case complexity is at least $\\Theta(\\log n)$ on finite $\\Delta$-regular trees, but their average local complexity is constant. \nSecond, we show that there are problems in the class $\\mathsf{BOREL}$ that are not in the class $\\mathsf{RLOCAL}(O(\\log\\log n)) = \\mathsf{RLOCAL}(o(\\log n))$, on $\\Delta$-regular trees. This shows that one cannot in general hope that results from (Borel) measurable combinatorics can be turned into very efficient (sublogarithmic) distributed algorithms. \n\n\\subsection{Preliminaries}\n\\label{subsec:preliminaries_separating}\n\nWe first show that there is a strong connection between randomized local complexities and uniform local complexities. Afterwards, we introduce a generic construction that turns any LCL into a new LCL. We later use this generic transformation to construct an LCL that is contained in the set $\\mathsf{BOREL} \\setminus \\mathsf{RLOCAL}(O(\\log \\log n))$.\n\n\\paragraph{Uniform vs Local Randomized Complexity} \n\nWe will now discuss the connections between uniform and randomized complexities. \nNote that the easy part of the connection between the two concepts is turning uniform local algorithms into randomized ones, as formalized in the following proposition. \n\n\\begin{proposition}\n\\label{prop:bee_gees}\nWe have $\\mathsf{OLOCAL}(t(\\varepsilon)) \\subseteq \\mathsf{RLOCAL}(t(1\/n^{O(1)}))$. \n\\end{proposition}\n\\begin{proof}\nWe claim that an uniform local algorithm $\\mathcal{A}$ with an uniform local complexity of $t(\\varepsilon)$ can be turned into a local randomized algorithm $\\mathcal{A}'$ with a local complexity of $t(1\/n^{O(1)})$. \n\nThe algorithm $\\mathcal{A}'$ simulates $\\mathcal{A}$ on an infinite $\\Delta$-regular tree -- each vertex $u$ of degree less than $\\Delta$ in the original tree pretends that the tree continues past its virtual half-edges and the random bits of $u$ are used to simulate the random bits in this virtual subtree. \nChoosing $\\varepsilon = 1\/n^{O(1)}$, one gets that the probability of $\\mathcal{A}'$ needing to look further than $t(\\varepsilon)$ for any vertex is bounded by $1\/n^{O(1)}$, as needed in the definition of the randomized local complexity. \n\\end{proof}\n\n\nOn the other hand, we will use the following proposition from \\cite{grebik_rozhon2021toasts_and_tails}. It informally states that the existence of \\emph{any} uniform local algorithm together with the existence of a sufficiently fast randomized local algorithm for a given LCL $\\Pi$ directly implies an upper bound on the uniform complexity of $\\Pi$.\n\\begin{proposition}[\\cite{grebik_rozhon2021toasts_and_tails}]\n\\label{prop:local->uniform}\nLet $\\mathcal{A}$ be an uniform local algorithm solving an LCL problem $\\Pi$ such that its \\emph{randomized} local complexity on finite $\\Delta$-regular trees is $t(n)$ for $t(n) = o(\\log n)$. \nThen, the uniform local complexity of $\\mathcal{A}$ on infinite $\\Delta$-regular trees is $O(t(1\/\\varepsilon))$.\n\\end{proposition}\n\n\\cref{prop:local->uniform} makes our life simpler, since we can take a known randomized distributed local algorithm with local complexity $g(n) = o(\\log n)$ and check if it works without the knowledge of $n$. If yes, this automatically implies that the uniform local complexity of the algorithm is $h(\\varepsilon) = O(g(1\/\\varepsilon))$. \nIn particular, combining \\cref{prop:local->uniform} with the work of \\cite{chang_kopelowitz_pettie2019exp_separation,chang_pettie2019time_hierarchy_trees_rand_speedup} and verifying that the algorithms of \nFischer and Ghaffari~\\cite[Section 3.1.1]{fischer-ghaffari-lll} and Chang~et~al.~\\cite[Section 5.1]{ChangHLPU20} work without the knowledge of $n$, we obtain the following result. \n\n\n\\begin{theorem}\n\\label{thm:classification_local_uniform_trees}\nWe have:\n\\begin{itemize}\n \\item $\\mathsf{LOCAL}(O(1)) = \\mathsf{OLOCAL}(O(1))$\n \\item $\\mathsf{RLOCAL}(O(\\log^* n)) = \\mathsf{OLOCAL}(O(\\log^* 1\/\\varepsilon))$\n \\item $\\mathsf{RLOCAL}(O(\\log\\log n)) = \\mathsf{OLOCAL}(O(\\log\\log 1\/\\varepsilon))$\n\\end{itemize}\nMoreover, there are no other possible uniform local complexities for $t(\\varepsilon) = o(\\log 1\/\\varepsilon)$. \n\\end{theorem}\nWe note that the first two items are proven in \\cite{grebik_rozhon2021toasts_and_tails}.\nFor the third item it suffices by known reductions to find an uniform algorithm\nsolving a version of the distributed Lov\\'asz Local Lemma (LLL) on so-called tree-structured dependency graphs considered in~\\cite{ChangHLPU20}.\n\n\\begin{proof}[Proof sketch]\nThe proof follows from \\cref{prop:bee_gees} and the following ideas. \nThe first item follows from the fact that any local algorithm with local complexity $O(1)$ can simply be made uniform. More specifically, there exists a constant $n_0$ --- depending only on the problem $\\Pi$ --- such that the algorithm, being told the size of the graph is $n_0$, is correct on \\emph{any} graph of size $n \\ge n_0$. \n\nFor the second item, by the work of \\cite{chang_kopelowitz_pettie2019exp_separation,chang_pettie2019time_hierarchy_trees_rand_speedup}, it suffices to check that there is an uniform distributed $(\\Delta + 1)$-coloring algorithm with uniform local complexity $O(\\log^* 1\/\\varepsilon)$. Such an algorithm was given in \\cite{HolroydSchrammWilson2017FinitaryColoring}, or follows from the work in \\cite{Korman_Sereni_Viennot2012Pruning_algorithms_+_oblivious_coloring} and \\cref{prop:local->uniform}. \n\nSimilarly, for the third item, by the work of \\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} it suffices to check that there is an uniform distributed algorithm for a specific LLL problem on trees\nwith uniform local complexity $O(\\log\\log 1\/\\varepsilon)$.\nSuch an algorithm can be obtained by combining the randomized \\emph{pre-shattering} algorithm of Fischer and Ghaffari~\\cite[Section 3.1.1]{fischer-ghaffari-lll} and the deterministic \\emph{post-shattering} algorithm of Chang~et~al.~\\cite[Section 5.1]{ChangHLPU20} in a \\emph{graph shattering} framework~\\cite{barenboim2016locality}, which solves the LLL problem \nwith local complexity $O(\\log\\log n)$.\nBy \\cref{prop:local->uniform}, it suffices to check that this algorithm can be made to work even if it does not know the size of the graph $n$. We defer the details to \\cref{sec:LLL}.\nThis finishes the proof of \\cref{thm:classification_local_uniform_trees}. \n\\end{proof}\n\n\n\n\n\n\n\n\\paragraph{Adding Paths}\nBefore we proceed to show some separation results, we define a certain construction that turns any LCL problem $\\Pi$ into a new LCL problem $\\overline{\\Pi}$ with the following property. If the original problem $\\Pi$ cannot be solved by a fast local algorithm, then the same holds for $\\overline{\\Pi}$. However, $\\overline{\\Pi}$ might be strictly easier to solve than $\\Pi$ for $\\mathsf{BOREL}$ constructions. \n\n\\begin{definition}\nLet $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ be an LCL. \nWe define an LCL $\\overline{\\Pi}=(\\Sigma',\\mathcal{V}',\\mathcal{E}')$ as follows.\nLet $\\Sigma'$ be $\\Sigma$ together with one new label.\nLet $\\mathcal{V}'$ be the union of $\\mathcal{V}$ together with any cardinality-$\\Delta$ multiset that contains the new label exactly two times.\nLet $\\mathcal{E}'$ be the union of $\\mathcal{E}$ together with the cardinality-$2$ multiset that contains the new label twice.\n\\end{definition}\n\nIn other words, the new label determines doubly infinite lines in infinite $\\Delta$-regular trees, or lines that start and end in virtual half-edges in finite $\\Delta$-regular trees.\nMoreover, a vertex that is on such a line does not have to satisfy any other vertex constraint.\nWe call these vertices \\emph{line-vertices} and each edge on a line a \\emph{line-edge}.\n\n\n\\begin{proposition}\n\\label{pr:AddingLineLLL}\nLet $\\Pi$ be an LCL problem such that $\\overline{\\Pi}\\in \\mathsf{RLOCAL}(t(n))$ for $t(n) = o(\\log(n))$. Then also $\\Pi \\in \\mathsf{RLOCAL}(O(t(n)))$.\n\\end{proposition}\n\\begin{proof}\nLet $\\overline{\\mathcal{A}}$ be a randomized $\\mathsf{LOCAL}$ algorithm that solves $\\overline{\\Pi}$ in $t(n) = o(\\log n)$ rounds with probability at least $1 - 1\/n^C$ for some sufficiently large constant $C$. We will construct an algorithm $\\mathcal{A}$ for $\\Pi$ with complexity $t(n)$ that is correct with probability $1 - \\frac{4}{n^{C\/3 - 2}}$. The success probability of $\\mathcal{A}$ can then be boosted by ``lying to it'' that the number of vertices is $n^{O(1)}$ instead of $n$; this increases the running time by at most a constant factor and boosts the success probability back to $1 - 1\/n^C$. \n\nConsider a $\\Delta$-regular rooted finite tree $T$ of depth $10t(n)$ and let $u$ be its root vertex. \nNote that $|T| \\le \\Delta^{10t(n)+1} < n$, for $n$ large enough. \n\n\nWe start by proving that when running $\\overline{\\mathcal{A}}$ on the tree $T$, then $u$ is most likely not a line-vertex. This observation then allows us to turn $\\overline{\\mathcal{A}}$ into an algorithm $\\mathcal{A}$ that solves $\\Pi$ on any $\\Delta$-regular input tree.\nLet $X$ be the indicator of $\\mathcal{A}$ marking $u$ as a line-vertex. Moreover, for $i \\in [\\Delta]$, let $Y_i$ be the indicator variable for the following event. The number of line edges in the $i$-th subtree of $u$ with one endpoint at depth $5t(n)$ and the other endpoint at depth $5t(n) + 1$ is odd.\n\nBy a simple parity argument we can relate the value of $X$ with the values of $Y_1,Y_2,\\ldots,Y_\\Delta$ as follows. If $X = 0$, that is, $u$ is not a line-vertex, then all of the $Y_i$'s have to be $0$, as each path in the tree $T$ is completely contained in one of the $\\Delta$ subtrees of $u$. On the other hand, if $u$ is a line-vertex, then there exists exactly one path, the one containing $u$, that is not completely contained in one of the $\\Delta$ subtrees of $u$. This in turn implies that exactly two of the $\\Delta$ variables $Y_1,Y_2,\\ldots,Y_\\Delta$ are equal to $1$.\n\nThe random variables $Y_1,Y_2,\\ldots,Y_\\Delta$ are identically distributed and mutually independent. Hence, if $\\P(Y_i = 1) > \\frac{1}{n^{C\/3}}$, then the probability that there are at least $3$ $Y_i$'s equal to $1$ is strictly greater than $\\left(\\frac{1}{n^{C\/3}}\\right)^3 = \\frac{1}{n^C}$. This is a contradiction, as in that case $\\overline{\\mathcal{A}}$ does not produce a valid output, which according to our assumption happens with probability at most $\\frac{1}{n^C}$. \n\nThus, we can conclude that $\\P(Y_i = 1) \\leq \\frac{1}{n^{C\/3}}$. By a union bound, this implies that all of the $Y_i$'s are zero with probability at least $\\frac{1}{n^{C\/3 - 1}}$, and in that case $u$ is not a line-vertex. \n\nFinally, the algorithm $\\mathcal{A}$ simulates $\\overline{\\mathcal{A}}$ as if the neighborhood of each vertex was a $\\Delta$-regular branching tree up to depth at least $t(n)$. \n\nIt remains to analyze the probability that $\\mathcal{A}$ produces a valid solution for the LCL problem $\\Pi$. To that end, let $v$ denote an arbitrary vertex of the input tree. The probability that the output of $v$'s half edges satisfy the vertex constraint of the LCL problem $\\overline{\\Pi}$ is at least $1 - 1\/n^C$. Moreover, in the case that $v$ is not a line-vertex, which happens with probability at least $1 - \\frac{1}{n^{C\/3 - 1}}$, the output of $v$'s half edges even satisfy the vertex constraint of the LCL problem $\\Pi$. Hence, by a union bound the vertex constraint of the LCL problem $\\Pi$ around $v$ is satisfied with probability at least $1 - \\frac{2}{n^{C\/3 - 1}}$. With exactly the same reasoning, one can argue that the probability that the output of $\\mathcal{A}$ satisfies the edge constraint of the LCL problem $\\Pi$ at a given edge is also at least $1 - \\frac{2}{n^{C\/3 - 1}}$. Finally, doing a union bound over the $n$ vertex constraints and $n-1$ edge constraints it follows that $\\mathcal{A}$ produces a valid solution for $\\Pi$ with probability at least $1 - \\frac{4}{n^{C\/3 - 2}}$.\n\\end{proof}\n\n\\begin{remark}\nThe ultimate way how to solve LCLs of the form $\\overline{\\Pi}$ is to find a spanning forest of lines.\nThis is because once we have a spanning forest of lines, then we might declare every vertex to be a line-vertex and label every half-edge that is contained on a line with the new symbol in $\\overline{\\Pi}$.\nThe remaining half-edges can be labeled arbitrarily in such a way that the edge constraints are satisfied.\n\nPut otherwise, once we are able to construct a spanning forest of lines and $\\mathcal{E}\\not=\\emptyset$, where $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$, then $\\overline{\\Pi}$ can be solved.\nMoreover, in that case the complexity of $\\overline{\\Pi}$ is upper bounded by the complexity of finding the spanning forest of lines.\n\nLyons and Nazarov \\cite{lyons2011perfect} showed that $\\Pi_{\\text{pm}}\\in \\mathsf{fiid}$ and it is discussed in \\cite{LyonsTrees} that the construction can be iterated $(\\Delta-2)$-many times. After removing each edge that is contained in one of the $\\Delta-2$ perfect matchings each vertex has a degree of two.\nHence, it is possible to construct a spanning forest of lines as fiid.\nConsequently, $\\overline{\\Pi}\\in \\mathsf{fiid}$ for every $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ with $\\mathcal{E}\\not=\\emptyset$.\n\\end{remark}\n\n\n\\subsection{$\\mathsf{LOCAL}(O(\\log^* n)) \\not= \\mathsf{BOREL}$}\n\nWe now formalize the proof that $\\mathsf{LOCAL}(O(\\log^* n)) \\not= \\mathsf{BOREL}$ from \\cref{sec:introseparation}. \nRecall that $\\Pi_{\\chi,\\Delta}$ is the proper vertex $\\Delta$-coloring problem.\nThe following claim directly follows by \\cref{pr:AddingLineLLL} together with the fact that $\\Pi_{\\chi,\\Delta}\\not\\in \\mathsf{LOCAL}(O(\\log^* n))$.\n\n\\begin{claim}\nWe have $\\overline{\\Pi_{\\chi,\\Delta}}\\not \\in \\mathsf{LOCAL}(O(\\log^* n))$.\n\\end{claim}\n\nOn the other hand we show that adding lines helps to find a Borel solution.\nThis already provides a simple example of a problem in the set $\\mathsf{BOREL} \\setminus \\mathsf{LOCAL}(O(\\log^* n))$.\n\n\\begin{proposition}\nWe have $\\overline{\\Pi_{\\chi,\\Delta}} \\in \\mathsf{BOREL}$. \n\\end{proposition}\n\\begin{proof}\nBy \\cite{KST}, every Borel graph of finite maximum degree admits a Borel maximal independent set.\nLet $\\mathcal{G}$ be a Borel $\\Delta$-regular acyclic graph on a standard Borel space $X$.\nBy an iterative application of the fact above we find Borel sets $A_1,\\dots,A_{\\Delta-2}$ such that $A_1$ is a maximal independent set in $\\mathcal{G}$ and $A_i$ is a maximal independent set in $\\mathcal{G}\\setminus \\bigcup_{j1$.\nWe think of $A_1,\\dots,A_{\\Delta-2}$ as color classes for the first $\\Delta-2$ colors.\nLet $B=X\\setminus \\bigcup_{i\\in [\\Delta-2]}A_i$ and let $\\mathcal{H}$ be the subgraph of $\\mathcal{G}$ determined by $B$.\nIt is easy to see that the maximum degree of $\\mathcal{H}$ is $2$.\nIn particular, $\\mathcal{H}$ consists of finite paths, one-ended infinite paths or doubly infinite paths.\nWe use the extra label in $\\overline{\\Pi_{\\chi,\\Delta}}$ to mark the doubly infinite paths.\nIt remains to use the remaining $2$ colors in $[\\Delta]$ to define a proper vertex $2$-coloring of the finite and one-ended paths.\nIt is a standard argument that this can be done in a Borel way and it is easy to verify that, altogether, we found a Borel $\\overline{\\Pi_{\\chi,\\Delta}}$-coloring of $\\mathcal{G}$.\n\\end{proof}\n\n\n\n\\subsection{Examples and Lower Bound}\n\nIn this subsection we define two LCLs and show that they are not in the class $\\mathsf{RLOCAL}(o(\\log n))$.\nIn the next subsection we show that one of them is in the class $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon)$ and the other in the class $\\mathsf{BOREL}$.\nBoth examples that we present are based on the following relaxation of the perfect matching problem.\n\n\\begin{definition}[Perfect matching in power-$2$ graph]\nLet $\\Pi^2_{\\text{pm}}$ be the perfect matching problem in the power-$2$ graph, i.e., in the graph that we obtain from the input graph by adding an edge between any two vertices that have distance at most $2$ in the input graph\n\\end{definition}\n\nWe show that $\\Pi^2_{\\text{pm}}$ is not contained in $\\mathsf{RLOCAL}(o(\\log n))$. \\cref{pr:AddingLineLLL} then directly implies that $\\overline{\\Pi^2_{\\text{pm}}}$ is not contained in $\\mathsf{RLOCAL}(o(\\log n))$ as well.\nOn the other hand, we later use a one-ended spanning forest decomposition to show that $\\Pi^2_{\\text{pm}}$ is in $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon)$ and a one or two-ended spanning forest decomposition to show that $\\overline{\\Pi^2_{\\text{pm}}}$ is in $\\mathsf{BOREL}$.\n\n\nWe first show the lower bound result.\nThe proof is based on a simple parity argument as in the proof of \\cref{pr:AddingLineLLL}.\n\n\\begin{theorem}\nThe problems $\\Pi^2_{\\text{pm}}$ and $\\overline{\\Pi^2_{\\text{pm}}}$ are not in the class $\\mathsf{RLOCAL}(o(\\log n))$.\n\\end{theorem}\n\n\\begin{proof}\nSuppose that the theorem statement does not hold for $\\Pi^2_{\\text{pm}}$.\nThen, by \\cref{thm:basicLOCAL}, there is a distributed randomized algorithm solving $\\Pi^2_{\\text{pm}}$ with probability at least $1 - 1\/n$. By telling the algorithm that the size of the input graph is $n^{2\\Delta}$ instead of $n$, we can further boost the success probability to $1 - 1\/n^{2\\Delta}$. The resulting round complexity is at most a constant factor larger, as $\\Delta = O(1)$.\nWe can furthermore assume that the resulting algorithm $\\mathcal{A}$ will always output for each vertex $v$ a vertex $M(v) \\neq v$ in $v$'s $2$-hop neighborhood, even if $\\mathcal{A}$ fails to produce a valid solution. The vertex $M(v)$ is the vertex $v$ decides to be matched to, and if $\\mathcal{A}$ produces a valid solution it holds that $M(M(v)) = v$.\n\nConsider a fully branching $\\Delta$-regular tree $T$ of depth $10t(n)$. Note that $|T| < n$ for $n$ large enough. Let $u$ be the root of $T$ and $T_i$ denote the $i$-th subtree of $u$ (so that $u \\not\\in T_i$).\nLet $S_i$ denote the set of vertices $v$ in $T_i$ such that both $v$ and $M(v)$ have distance at most $2t(n)$ from $u$.\nNote that $M(v)$ is not necessarily a vertex of $T_i$.\nLet $Y_i$ be the indicator of whether $S_i$ is even.\n\nObserve that, if $\\mathcal{A}$ does not fail on the given input graph, then the definition of the $S_i$ implies that every vertex in $\\{ u \\} \\cup \\bigcup_{i \\in [\\Delta]} S_i$ is matched to a vertex in $\\{ u \\} \\cup \\bigcup_{i \\in [\\Delta]} S_i$.\nHence, the number of vertices in $\\{ u \\} \\cup \\bigcup_{i \\in [\\Delta]} S_i$ is even, which implies that we cannot have $Y_1 = Y_2 = \\ldots = Y_\\Delta = 1$ unless $\\mathcal{A}$ fails.\nNote that $Y_i$ depends only on the output of $\\mathcal{A}$ at the vertices in $T_i$ that have a distance of precisely $2t(n) - 1$ or $2t(n)$ to $u$ (as all other vertices in $T_i$ are guaranteed to be in $S_i$). \nHence, the events $Y_i = 1$, $i \\in [\\Delta]$, are independent, and, since $\\P(Y_1 = 1) = \\ldots = \\P(Y_{\\Delta} = 1)$ (as all vertices in all $T_i$ see the same topology in their $t(n)$-hop view), we obtain $(\\P(Y_1 = 1))^{\\Delta} \\le 1\/n^{2\\Delta}$, which implies $\\P(Y_i = 1) \\le 1\/n^2$, for any $i \\in [\\Delta]$.\n\n\nHence, with probability at least $1 - \\Delta \/n^2$ we have $Y_1=Y_2=\\ldots = Y_\\Delta = 0$, by a union bound.\nLet $v_1, v_2, \\dots, v_\\Delta$ denote the neighbors of $u$ with $v_i$ being the root of subtree $T_i$.\nNote that if $\\mathcal{A}$ does not fail and $Y_1=Y_2=\\ldots = Y_\\Delta = 0$, then it has to be the case that $u$ is matched to a vertex in some subtree $T_i$ (not necessarily $v_i$), while any vertex from $N_{-i}(u) := \\{v_1, v_2, \\dots, v_{i-1}, v_{i+1}, \\dots, v_\\Delta\\}$ is matched with another vertex from $N_{-i}(u)$ (hence, we already get that $\\Delta$ needs to be odd).\nThis is because each subtree needs to have at least one vertex matched to a different subtree (since $|S_i|$ is odd for each $i \\in [\\Delta]$).\nIf $M(u)$ is a vertex in $T_i$ and, for any $v \\in N_{-i}(u)$, we have $M(v) \\in N_{-i}(u)$, we say that $u$ orients the edge going to $v_i$ outwards.\n\nConsider a path $(u_0 = u, u_1, \\dots, u_k)$, where $k = 2t(n) + 3$.\nBy the argument provided above, the probability that $u$ orients an edge outwards is at least $1 - \\Delta\/n^2 - 1\/n^{2\\Delta}$, and, if $u$ orients an edge outwards, the probability that this edge is not $uu_1$ is $(\\Delta - 1)\/\\Delta$, by symmetry.\nHence, we obtain (for sufficiently large $n$) that $\\P(\\mathcal{E}_1) \\geq 1\/2$, where $\\mathcal{E}_1$ denotes the event that vertex $u$ orients an edge different from $uu_1$ outwards.\nWith an analogous argument, we obtain that $\\P(\\mathcal{E}_2) \\geq 1\/2$, where $\\mathcal{E}_2$ denotes the event that vertex $u_k$ orients an edge different from $u_ku_{k-1}$ outwards.\nSince $\\mathcal{E}_1$ and $\\mathcal{E}_2$ are independent (due to the fact that the distance between any neighbor of $u$ and any neighbor of $u_k$ is at least $2t(n) + 1$), we obtain $\\P(\\mathcal{E}_1 \\cap \\mathcal{E}_2) \\geq 1\/4$.\nMoreover, the probability that all vertices in $\\{ u_1, \\dots, u_{k-1} \\}$ orient an edge outwards is at least $1 - (2t(n)+2)(\\Delta\/n^2 + 1\/n^{2\\Delta})$, which is at least $4\/5$ for sufficiently large $n$.\nThus, by a union bound, there is a probability of at least $1 - 3\/4 - 1\/5 - 1\/n^{2\\Delta}$ that $\\mathcal{A}$ does not fail, all vertices in $(u, u_1, \\dots, u_k)$ orient an edge outwards, and the outward oriented edges chosen by $u$ and $u_k$ are not $uu_1$ and $u_ku_{k-1}$, respectively.\nFor sufficiently large $n$, the indicated probability is strictly larger than $0$.\nWe will obtain a contradiction and conclude the proof for $\\Pi^2_{\\text{pm}}$ by showing that there is no correct solution with the described properties.\n\nAssume in the following that the solution provided by $\\mathcal{A}$ is correct and $u, u_1, \\dots, u_k$ satisfy the mentioned properties.\nSince $u$ orients an edge different from $uu_1$ outwards, $u_1$ must be matched to a neighbor of $u$, which implies that $u_1$ orients edge $u_1u$ outwards.\nUsing an analogous argumentation, we obtain inductively that $u_j$ orients edge $u_ju_{j-1}$ outwards, for all $2 \\leq j \\leq k$.\nHowever, this yields a contradiction to the fact that the edge that $u_k$ orients outwards is different from $u_ku_{k-1}$, concluding the proof (for $\\Pi^2_{\\text{pm}}$).\n\nThe theorem statement for $\\overline{\\Pi^2_{\\text{pm}}}$ follows by applying \\cref{pr:AddingLineLLL}.\n\n\n\\end{proof}\n\n\n\n\n\n\\subsection{Upper Bounds Using Forest Decompositions}\n\nWe prove an upper bound for the problems $\\Pi^2_{\\text{pm}}$ and $\\overline{\\Pi^2_{\\text{pm}}}$ defined in the previous subsection.\nWe use a technique of decomposing the input graph in a spanning forest with some additional properties, i.e., one or two ended. This technique was used in \\cite{BrooksMeas} to prove Brooks' theorem in a measurable context.\nNamely, they proved that if $\\mathcal{G}$ is a Borel $\\Delta$-regular acyclic graph and $\\mu$ is a Borel probability measure, then it is possible to erase some edges of $\\mathcal{G}$ in such a way that the remaining graph is a one-ended spanning forest on a $\\mu$-conull set.\nIt is not hard to see that one can solve $\\Pi^2_{\\text{pm}}$ on one-ended trees in an inductive manner, starting with the leaf vertices. Consequently we have $\\Pi^2_{\\text{pm}}\\in \\mathsf{MEASURE}$.\nWe provide a quantitative version of this result and thereby show that the problem of constructing a one-ended forest is contained in $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon)$.\nA variation of the construction shows that it is possible to find a one or two-ended spanning forest in a Borel way.\nThis allows us to show that $\\overline{\\Pi^2_{\\text{pm}}}\\in \\mathsf{BOREL}$.\nAs an application of the decomposition technique we show that Vizing's theorem, i.e., proper $\\Delta+1$ edge coloring $\\Pi_{\\chi',\\Delta+1}$, for $\\Delta=3$ is in the class $\\mathsf{OLOCAL}(\\operatorname{poly} \\log 1\/\\varepsilon)$.\n\n\n\\subsubsection{Uniform Complexity of the One-Forest Decomposition }\nWe start with the definition of a one-ended spanning forest.\nIt can be viewed as a variation of the edge grabbing problem $\\Pi_{\\text{edgegrab}}$.\nNamely, if a vertex $v$ grabs an incident edge $e$, then we think of $e$ as being oriented away from $v$ and $v$ selecting the other endpoint of $e$ as its parent.\nSuppose that $\\mathcal{T}$ is a solution of $\\Pi_{\\text{edgegrab}}$ on $T_\\Delta$.\nWe denote with $\\mathcal{T}^{\\leftarrow}(v)$ the subtree with root $v$.\n\n\n\\begin{definition}[One-ended spanning forest]\nWe say that a solution $\\mathcal{T}$ of $\\Pi_{\\text{edgegrab}}$ is a \\emph{one-ended spanning forest} if $\\mathcal{T}^{\\leftarrow}(v)$ is finite for every $v\\in T_\\Delta$ (see \\cref{fig:one_ended}). \n\n\\end{definition}\n\n\\begin{figure}\n \\centering\n \\includegraphics[width=.4\\textwidth]{fig\/one_ended.pdf}\n \\caption{An example of a one-ended forest decomposition of an infinite 3-regular tree. }\n \\label{fig:one_ended}\n\\end{figure}\n\n\nNote that on an infinite $\\Delta$-regular tree every connected component of a one-ended spanning forest must be infinite.\nFurthermore, recall that we discussed above that it is possible to construct a one-ended spanning forest in $\\mathsf{MEASURE}$ by \\cite{BrooksMeas}.\nNext we formulate our quantitative version.\n\n\n\n\n\n\\begin{restatable}{theorem}{oneendedtreetail}\n\\label{thm:one_ended_forest}\nThe one-ended spanning forest can be constructed by an uniform local algorithm with an uniform local complexity of $O(\\operatorname{poly}\\log1\/\\varepsilon)$. \nMore precisely, there is an uniform distributed algorithm $\\mathcal{A}$ that computes a one-ended spanning forest and if we define $R(v)$ to be the smallest coding radius that allows $v$ to compute $\\mathcal{T}^{\\leftarrow}(v)$, then\n\\[\nP(R(v) > O(\\operatorname{poly} \\log 1\/\\varepsilon)) \\le \\varepsilon. \n\\]\n\\end{restatable}\n\nThe formal proof of \\cref{thm:one_ended_forest} can be found in \\cref{sec:formalproof}.\nWe first show the main applications of the result.\nNamely, we show that $\\Pi^2_{\\text{pm}}$ can be solved inductively starting from the leaves of the one-ended trees. This proves that $\\Pi^2_{\\text{pm}}$ is in $\\mathsf{OLOCAL}(\\operatorname{poly} \\log 1\/\\varepsilon)$ and hence $\\mathsf{OLOCAL}(\\operatorname{poly} \\log 1\/\\varepsilon) \\setminus \\mathsf{RLOCAL}(o(\\log n))$ is non-empty.\n\n\\begin{theorem}\\label{thm:twomatchFFIID}\nThe problem $\\Pi^2_{\\text{pm}}$ is in the class $\\mathsf{MEASURE}$ and $\\mathsf{OLOCAL}(O(\\operatorname{poly} \\log 1\/\\varepsilon))$.\n\\end{theorem}\n\\begin{proof}\nFirst we show that every infinite one-ended directed tree $T$ admits an inductive solution from the leaves to the direction of infinity. More precisely, we define a power-$2$ perfect matching on the infinite one-ended directed tree $T$ such that the following holds. Let $v$ be an arbitrary vertex and $T^{\\leftarrow}(v)$ the finite subtree rooted at $v$. Then, all vertices in $T^{\\leftarrow}(v)$, with the possible exception of $v$, are matched to a vertex in $T^{\\leftarrow}(v)$. Moreover, for each vertex in $T^{\\leftarrow}(v) \\setminus \\{v\\}$, it is possible to determine to which vertex it is matched by only considering the subtree $T^{\\leftarrow}(v)$. \nWe show that it is possible to define a perfect matching in that way by induction on the height of the subtree $T^{\\leftarrow}(v)$. \n\nWe first start with the leaves and declare them as unmatched in the first step of the induction.\n\nNow, consider an arbitrary vertex $v$ with $R \\geq 1$ children that we denote by $v_1, v_2, \\ldots, v_R$.\nBy the induction hypothesis, all vertices in $T^{\\leftarrow}(v_i) \\setminus \\{v_i\\}$ are matched with a vertex in $T^{\\leftarrow}(v_i)$ for $i \\in [R]$ and it is possible to compute that matching given $T^{\\leftarrow}(v)$. However, some of the children of $v$ are possibly still unmatched. If the number of unmatched children of $v$ is even, then we simply pair them up in an arbitrary but fixed way and we leave $v$ unmatched. Otherwise, if the number of unmatched children of $v$ is odd, we pair up $v$ with one of its unmatched children and we pair up all the remaining unmatched children of $v$ in an arbitrary but fixed way. This construction clearly satisfies the criteria stated above. In particular, the resulting matching is a perfect matching, as every vertex gets matched eventually.\n\n\nBy \\cref{thm:one_ended_forest} there is an uniform distributed algorithm of complexity $O(\\operatorname{poly} \\log 1\/\\varepsilon)$ that computes a one-ended spanning forest $\\mathcal{T}$ on $T_\\Delta$.\nMoreover, every vertex $v\\in T_\\Delta$ can compute where it is matched with an uniform local complexity of $O(\\operatorname{poly} \\log 1\/\\varepsilon)$.\nThis follows from the discussion above, as a vertex $v$ can determine where it is matched once it has computed $\\mathcal{T}^{\\leftarrow}(w)$ for each $w$ in its neighborhood.\nHence, $\\Pi^2_{\\text{pm}}$ is in the class $\\mathsf{OLOCAL}(O(\\operatorname{poly} \\log 1\/\\varepsilon))$.\n\nSimilarly, the one-ended spanning forest $\\mathcal{T}\\subseteq \\mathcal{G}$ can be constructed in the class $\\mathsf{MEASURE}$ by \\cite{BrooksMeas}.\nBy the discussion above, this immediately implies that $\\Pi^2_{\\text{pm}}$ is in the class $\\mathsf{MEASURE}$. \n\\end{proof}\n\n\n\n\n\\subsubsection{Decomposition in $\\mathsf{BOREL}$}\nNext, we show that a similar, but weaker, decomposition can be done in a Borel way.\nNamely, we say that a subset of edges of an acyclic infinite graph $G$, denoted by $\\mathcal{T}$, is a \\emph{one or two-ended spanning forest} if every vertex $v\\in G$ is contained in an infinite connected component of $\\mathcal{T}$ and each infinite connected component of $\\mathcal{T}$ has exactly one or two directions to infinity.\nLet $S$ be such a connected component.\nNote that if $S$ has one end, then we can solve $\\Pi^2_{\\text{pm}}$ on $S$ as above.\nIf $S$ has two ends, then $S$ contains a doubly infinite path $P$ with the property that erasing $P$ splits $S\\setminus P$ into finite connected components.\nIn order to solve $\\overline{\\Pi^2_{\\text{pm}}}$ we simply solve $\\Pi^2_{\\text{pm}}$ on the finite connected components of $S\\setminus P$ (with some of the vertices in $S\\setminus P$ being potentially matched with vertices on the path $P$) and declare vertices on $P$ to be line vertices.\n\nThe high-level idea to find a one or two-ended spanning forest is to do the same construction as in the measure case and understand what happens with edges that do not disappear after countably many steps.\nA slight modification in the construction guarantees that what remains are doubly infinite paths.\n\n\\begin{restatable}{theorem}{oneendedforestborel}\n\\label{thm:OneTwoEndBorel}\nLet $\\mathcal{G}$ be a Borel $\\Delta$-regular forest. \nThen there is a Borel one or two-ended spanning forest $\\mathcal{T}\\subseteq \\mathcal{G}$.\n\\end{restatable}\n\nThe proof of \\cref{thm:OneTwoEndBorel} can be found in \\cref{sec:oneortwo}.\nWe remark that the notation used in the proof is close to the notation in the proof that gives a measurable construction of a one-ended spanning forest in \\cite{BrooksMeas}.\nNext, we show more formally how \\cref{thm:OneTwoEndBorel} implies that $\\overline{\\Pi^2_{\\text{pm}}}\\in \\mathsf{BOREL}$.\n\n\\begin{theorem}\nThe problem $\\overline{\\Pi^2_{\\text{pm}}}$ is in the class $\\mathsf{BOREL}$.\n\\end{theorem}\n\\begin{proof}\n\nLet $\\mathcal{T}$ be the one or two-ended spanning forest given by \\cref{thm:OneTwoEndBorel} and $S$ be a connected component of $\\mathcal{T}$.\nIf $S$ is one-ended, then we use the first part of the proof of \\cref{thm:twomatchFFIID} to find a solution of $\\Pi^2_{\\text{pm}}$ on $S$.\nSince deciding that $S$ is one-ended as well as computing an orientation towards infinity can be done in a Borel way, this yields a Borel labeling.\n\nSuppose that $S$ is two-ended.\nThen, there exists a doubly infinite path $P$ in $S$ with the property that the connected components of $S\\setminus P$ are finite.\nMoreover, it is possible to detect $P$ in a Borel way.\nThat is, declaring the vertices on $P$ to be line vertices and using the special symbol in $\\overline{\\Pi^2_{\\text{pm}}}$ for the half-edges on $P$ yields a Borel measurable labeling.\nLet $C\\not=\\emptyset$ be one of the finite components in $S\\setminus P$ and $v\\in C$ be the vertex of distance $1$ from $P$.\nOrient the edges in $C$ towards $v$, this can be done in a Borel way since $v$ is uniquely determined for $C$.\nThen the first part of the proof of \\cref{thm:twomatchFFIID} shows that one can inductively find a solution to $\\Pi^2_{\\text{pm}}$ on $C$ in such a way that all vertices, possibly up to $v$, are matched.\nIf $v$ is matched, then we are done.\nIf $v$ is not matched, then we add the edge that connects $v$ with $P$ to the matching. Note that it is possible that multiple vertices are matched with the same path vertex, but this is not a problem according to the definition of $\\overline{\\Pi^2_{\\text{pm}}}$. \nIt follows that this defines a Borel function on $H(\\mathcal{G})$ that solves $\\overline{\\Pi^2_{\\text{pm}}}$.\n\\end{proof}\n\n\n\n\n\n\\subsubsection{Vizing's Theorem for $\\Delta=3$}\nFinding a measurable or local version of Vizing's Theorem, i.e., proper edge $(\\Delta+1)$-coloring $\\Pi_{\\chi',\\Delta+1}$, was studied recently in \\cite{grebik2020measurable,weilacher2021borel,BernshteynVizing}.\nIt is however not known, even on trees, whether $\\Pi_{\\chi',\\Delta+1}$ is in $\\mathsf{RLOCAL}(O(\\log\\log n))$.\nHere we use the one-ended forest construction to show that $\\Pi_{\\chi',\\Delta+1}$ is in the class $\\mathsf{OLOCAL}(\\operatorname{poly}\\log 1\/\\varepsilon)$ for $\\Delta=3$.\n\n\n\n\\iffalse\n\\paragraph{Composition of Uniform Local Algorithms}\nWhen designing algorithms, one often composes several simple algorithms to obtain a more sophisticated algorithm. For example, an algorithm $\\mathcal{A}$ might first invoke algorithm $\\mathcal{A}_1$ and after $\\mathcal{A}_1$ finished, $\\mathcal{A}$ invokes algorithm $\\mathcal{A}_2$ with the input of $\\mathcal{A}_2$ being equal to the output of $\\mathcal{A}_1$. Then, the final output of $\\mathcal{A}$ is simply the output of $\\mathcal{A}_2$. We call such a composition a sequential composition. In the context of the classic deterministic local model, the round complexity of $\\mathcal{A}$ is simply the sum of the round complexities of the two algorithms $\\mathcal{A}_1$ and $\\mathcal{A}_2$ that are used as a subroutine. Here, we prove two such composition results in the context of uniform local algorithms. These general results will help us in \\cref{app:one_ended} to analyze the uniform complexity of an uniform algorithm that computes a so-called one-ended forest in a principled manner.\n\n\n\n\\begin{lemma}[Parallel Composition]\n\\label{lem:parallel_composition}\nLet $A_1$ and $A_2$ be two distributed uniform algorithms with an uniform local complexity of $t_{A_1}(\\varepsilon)$ and $t_{A_2}(\\varepsilon)$, respectively. \nLet $A$ be the parallel composition of $A_1$ and $A_2$. That is, the coding radius of algorithm $\\mathcal{A}$ at a vertex $u$ is the smallest radius such that $A$ knows the local answer of both $A_1$ and $A_2$ at $u$. \nThen, $t_A(\\varepsilon) \\leq \\max(t_{A_1}(\\varepsilon\/2), t_{A_2}(\\varepsilon\/2))$.\n\\end{lemma}\n\\begin{proof}\nThe probability that the coding radius of algorithm $\\mathcal{A}_1$ at vertex $u$ is larger than $\\max(t_{A_1}(\\varepsilon\/2), t_{A_2}(\\varepsilon\/2))$ is at most $\\varepsilon\/2$ and the same holds for the coding radius of algorithm $\\mathcal{A}_2$ at vertex $u$. Hence, by a simple union bound, the probability that the coding radius of algorithm $\\mathcal{A}$ at a vertex $u$ is larger than $\\max(t_{A_1}(\\varepsilon\/2), t_{A_2}(\\varepsilon\/2))$ is at most $\\varepsilon\/2 + \\varepsilon\/2 = \\varepsilon$. \\\\ Hence, $t_A(\\varepsilon) \\leq \\max(t_{A_1}(\\varepsilon\/2), t_{A_2}(\\varepsilon\/2))$, as desired.\n\\end{proof}\n\n\\fi\n\n\n\n\\begin{proposition}\nLet $\\Delta=3$. We have $\\Pi_{\\chi',\\Delta+1} \\in\\mathsf{OLOCAL}(\\operatorname{poly} \\log 1\/\\varepsilon)$.\n\\end{proposition}\n\\begin{proof}[Proof sketch.]\nBy \\cref{thm:one_ended_forest}, we can compute a one-ended forest decomposition $\\mathcal{T}$ with uniform local complexity $O(\\operatorname{poly} \\log 1\/\\varepsilon)$. \nNote that every vertex has at least one edge in $\\mathcal{T}$, hence the edges in $G \\setminus \\mathcal{T}$ form paths. \nThese paths can be $3$-edge colored by using the uniform version of Linial's coloring algorithm. This algorithm has an uniform local complexity of $O(\\log^* 1\/\\varepsilon)$. By \\cref{lem:sequential_composition}, the overall complexity is $O(\\operatorname{poly} \\log (\\Delta^{\\log^* 1\/\\varepsilon} \/ \\varepsilon) + \\log^* 1\/\\varepsilon)) = O(\\operatorname{poly} \\log 1\/\\varepsilon)$. \n\nFinally, we color the edges of $T$. We start from the leaves and color the edges inductively. In particular, whenever we consider a vertex $v$ we color the at most two edges in $T^{\\leftarrow}(v)$ below $v$. \nNote that there always is at least one uncolored edge going from $v$ in the direction of $T$. Hence, we can color the at most two edges below $v$ in $T^{\\leftarrow}(v)$ greedily -- each one neighbors with at most $3$ colored edges at any time. \n\\end{proof}\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\subsection{Proof of \\cref{thm:one_ended_forest}} \n\\label{sec:formalproof}\n\n\nIn this section we formally prove \\cref{thm:one_ended_forest}. \n\n\n\\oneendedtreetail*\n\nWe follow the construction of the one-ended spanning forest in \\cite{BrooksMeas}.\nThe construction proceeds in rounds. Before each round, some subset of the vertices have already decided which incident edge to grab, and we refer to these vertices as settled vertices. All the remaining vertices are called unsettled. The goal in each round is to make a large fraction of the unsettled vertices settled. To achieve this goal, our construction relies on certain properties that the graph induced by all the unsettled vertices satisfies. \n\nOne important such property is that the graph induced by all the unsettled vertices is expanding. That is, the number of vertices contained in the neighborhood around a given vertex grows exponentially with the radius. The intuitive reason why this is a desirable property is the following. Our algorithm will select a subset of the unsettled vertices and clusters each unsettled vertex to the closest selected vertex. As the graph is expanding and any two vertices in the selected subset are sufficiently far away, there will be a lot of edges leaving a given cluster. For most of these inter-cluster edges, both of the endpoints will become settled. This in turn allows one to give a lower bound on the fraction of vertices that become settled in each cluster.\n\nTo ensure that the graph induced by all the unsettled vertices is expanding, a first condition we impose on the graph induced by all the unsettled vertices is that it has a minimum degree of at least $2$. While this condition is a first step in the right direction, it does not completely suffice to ensure the desired expansion, as an infinite path has a minimum degree of $2$ but does not expand sufficiently. Hence, our algorithm will keep track of a special subset of the unsettled vertices, the so-called hub vertices. Each hub vertex has a degree of $\\Delta$ in the graph induced by all the unsettled vertices. That is, all of the neighbors of a hub vertex are unsettled as well. Moreover, each unsettled vertex has a bounded distance to the closest hub vertex, where the specific upper bound on the distance to the closest hub vertex increases with each round. As we assume $\\Delta > 2$, the conditions stated above suffice to show that the graph induced by all the unsettled vertices expands.\n\nNext, we explain in more detail how a vertex decides which edge to grab. Concretely, if a vertex becomes settled, it grabs an incident edge such that the other endpoint of that edge is strictly closer to the closest unsettled vertex as the vertex itself. For example, if a vertex becomes settled and there is exactly one neighbor that is still unsettled, then the vertex will grab the edge that the vertex shares with its unsettled neighbor. Grabbing edges in that way, one can show two things. First, for a settled vertex $v$, the set $T^{\\leftarrow}(v)$ of vertices behind $v$ does not change once $v$ becomes settled. The intuitive reason for this is that the directed edges point towards the unsettled vertices and therefore the unsettled vertices lie before $v$ and not after $v$. Moreover, one can also show that $T^{\\leftarrow}(v)$ only contains finitely many vertices. The reason for this is that at the moment a vertex $v$ becomes settled, there exists an unsettled vertex that is sufficiently close to $v$, where the exact upper bound on that distance again depends on the specific round in which $v$ becomes settled. \n\nWhat remains to be discussed is at which moment a vertex decides to become settled. As written above, in each round the algorithm considers a subset of the unsettled vertices and each unsettled vertex is clustered to the closest vertex in that subset. This subset only contains hub vertices and it corresponds to an MIS on the graph with the vertex set being equal to the set of hub vertices and where two hub vertices are connected by an edge if they are sufficiently close in the graph induced by all the unsettled vertices. Now, each cluster center decides to connect to exactly $\\Delta$ different neighboring clusters, one in each of its $\\Delta$ subtrees. Remember that as each cluster center is a hub vertex, all of its neighbors are unsettled as well. Now, each unsettled vertex becomes settled except for those that lie on the unique path between two cluster centers such that one of the cluster centers decided to connect to the other one.\n\n\n\n\n\n\n\\paragraph{Construction}\nThe algorithm proceeds in rounds. After each round, a vertex is either settled or unsettled and a settled vertex remains settled in subsequent rounds. Moreover, some unsettled vertices are so-called hub vertices. We denote the set of unsettled, settled and hub vertices after the $i$-th round with $U_i$, $S_i$ and $H_i$, respectively. We set $U_0 = H_0 = V$ prior to the first round and we always set $S_i = V \\setminus U_i$. \\cref{fig:constr1} illustrates the situation after round $i$. Moreover, we denote with $O_i$ a partial orientation of the edges of the infinite $\\Delta$-regular input tree $T$. In the beginning, $O_0$ corresponds to the partial orientation with no oriented edges. Each vertex is incident to at most $1$ outwards oriented edge in $O_i$. If a vertex is incident to an outwards oriented edge in $O_i$, then the other endpoint of the edge will be its parent in the one-ended forest decomposition. For each vertex $v$ in $S_i$, we denote with $T_{v,i}$ the smallest set that satisfies the following conditions. First, $v \\in T_{v,i}$. Moreover, if $u \\in T_{v,i}$ and $\\{w,u\\}$ is an edge that according to $O_i$ is oriented from $w$ to $u$, then $w \\in T_{v,i}$. We later show that $T_{v,i}$ contains exactly those vertices that are contained in the subtree $\\mathcal{T}^\\leftarrow(v)$.\n\n\n\\begin{figure}[H]\n \\centering\n \\includegraphics{fig\/one_ended_construction.pdf}\n \\caption{The picture illustrates the situation after some round $i$ of the construction. The vertices on the paths are all contained in $U_i$, while the vertices corresponding to the big black dots are additionally contained in $H_i$. Note that the length of each path is at least $d_i$. However, the distance between $u$ and $v$ in the original graph can be much smaller.}\n \\label{fig:constr1}\n\\end{figure}\n\n\nAfter the $i$-th round, the construction satisfies the following invariants.\n\n\n\n\\begin{enumerate}\n \n \\item For each $v \\in S_{i-1}$, we have $T_{v,i} = T_{v,i-1}$.\n \n \\item Let $v$ be an arbitrary vertex in $S_i$. Then, $T_{v,i}$ contains finitely many vertices and furthermore $T_{v,i} \\subseteq S_i$.\n \n \\item The minimum degree of the graph $T[U_i]$ is at least $2$ and each vertex in $H_i$ has a degree of $\\Delta$ in $T[U_i]$.\n \n \\item Each vertex in $U_i$ has a distance of at most $\\sum_{j=0}^i d_j$ to the closest vertex in $H_i$ in the graph $T[U_i]$.\n\\end{enumerate}\n\nWe now describe how to compute $U_i, S_i$, $H_i$ and $O_i$. The construction during the $i$-th round is illustrated in $\\cref{fig:constr2}$. Note that we can assume that the invariants stated above are satisfied after the $(i-1)$-th round. In the $i$-th round we have a parameter $d_i$ that we later set to $2^{2^i}$.\n\n\\begin{enumerate}\n \\item $H_i$ is a subset of $H_{i-1}$ that satisfies the following property. No two vertices in $H_i$ have a distance of at most $d_i$ in $T[U_{i-1}]$. Moreover, each vertex in $H_{i-1} \\setminus H_i$ has a distance of at most $d_i$ to the closest vertex in $H_i$ in the graph $T[U_{i-1}]$. Note that we can compute $H_i$ by computing an MIS in the graph with vertex set $H_{i-1}$ and where two vertices are connected iff they have a distance of at most $d_i$ in the graph $T[U_{i-1}]$. Hence, we can use Ghaffari's MIS algorithm \\cite{ghaffari2016MIS} to compute the set $H_i$.\n \\item Next, we describe how to compute $U_i$. \n We assign each vertex $u \\in U_{i-1}$ to the closest vertex in the graph $T[U_{i-1}]$ that is contained in $H_i$, with ties being broken arbitrarily. We note that there exists a node in $H_i$ with a distance of at most $(\\sum_{j=0}^{i-1} d_j) + d_i$ to $u$. To see why, note that Invariant $(4)$ from round $i - 1$ implies that there exists a vertex $w$ in $H_{i-1}$ with a distance of at most $\\sum_{j=0}^{i-1} d_j$ to $u$. We are done if $w$ is also contained in $H_i$. If not, then it follows from the way we compute $H_i$ that there exists a vertex in $H_i$ with a distance of at most $d_i$ to $w$, but then the triangle inequality implies that the distance from $u$ to that vertex is at most $(\\sum_{j=0}^{i-1} d_j) + d_i$, as desired. For each vertex $v \\in H_i$, we denote with $C(v)$ the set of all vertices in $U_{i-1}$ that got assigned to $v$. Now, let $E_v$ denote the set of edges that have exactly one endpoint in $C(v)$. We can partition $E_v$ into $E_{v,1} \\sqcup E_{v,2} \\sqcup \\ldots \\sqcup E_{v,\\Delta}$, where for $\\ell \\in [\\Delta]$, $E_{v,\\ell}$ contains all the edges in $E_v$ that are contained in the $\\ell$-th subtree of $v$. \n Invariants $(3)$ and $(4)$ after the $(i-1)$-th round imply that $E_{v,\\ell} \\neq \\emptyset$. Moreover, $E_{v,\\ell}$ contains only finitely many edges, as we have shown above that the cluster radius is upper bounded by $\\sum_{j=0}^{i} d_j$.\n \n Now, for $\\ell \\in [\\Delta]$, we choose an edge $e_\\ell$ uniformly at random from the set $E_{v,\\ell}$. Let $u_\\ell \\neq v$ be the unique vertex in $H_i$ such that one endpoint of $e_\\ell$ is contained in $C(u_\\ell)$. We denote with $P_{v,\\ell}$ the set of vertices that are contained in the unique path between $v$ and $u_\\ell$. Finally, we set $U_i = \\bigcup_{v \\in H_i, \\ell \\in [\\Delta]} P_{v,\\ell}$.\n \\item It remains to describe how to compute the partial orientation $O_i$. All edges that are oriented in $O_{i-1}$ will be oriented in the same direction in $O_i$. Additionally, we orient for each vertex $u$ that got settled in the $i$-th round, i.e., $u \\in S_i \\cap U_{i-1}$, exactly one incident edge away from $u$. Let $w$ be the closest vertex to $u$ in $T[U_{i-1}]$ that is contained in $U_i$, with ties being broken arbitrarily. We note that it follows from the discussions above that the distance between $u$ and $w$ in the graph $T[U_{i-1}]$ is at most $\\sum_{j=0}^{i} d_j$. We now orient the edge incident to $u$ that is on the unique path between $u$ and $w$ outwards. We note that this orientation is well-defined in the sense that we only orient edges that were not oriented before and that we don't have an edge such that both endpoints of that edge want to orient the edge outwards.\n\\end{enumerate}\n\n\n\\begin{figure}\n \\centering\n \\includegraphics{fig\/one_ended_construction2.pdf}\n \\caption{The picture illustrates the situation during some round $i$ of the construction. The vertices on all paths are contained in $U_{i-1}$ and the vertices corresponding to the big black dots are contained in $H_{i-1}$. The vertices $v$, $u_1$, $u_2$ and $u_3$ are the only vertices in the picture that are contained in $H_i$. All the vertices on the red paths remain unsettled during round $i$ while all the vertices on the black paths become settled. Each vertex that becomes settled during the $i$-th round orients one incident edge outwards, illustrated by the green edges. Each green edge is oriented towards the closest vertex that is still unsettled after the $i$-th round. We note that some of the edges in the picture are oriented away from $v$.}\n \\label{fig:constr2}\n\\end{figure}\n\n\nWe now prove by induction that our procedure satisfies all the invariants.\nPrior to the first round, all invariants are trivially satisfied. Hence, it remains to show that the invariants hold after the $i$-th round, given that the invariants hold after the $(i-1)$-th round. \n\n\\begin{enumerate}\n\n \\item Let $v \\in S_{i-1}$ be arbitrary. As $O_i$ is an extension of $O_{i-1}$, it follows from the definition of $T_{v,i-1}$ and $T_{v,i}$ that $T_{v,i-1} \\subseteq T_{v,i}$. Now assume that $T_{v,i} \\not \\subseteq T_{v,i-1}$. This would imply the existence of an edge $\\{u,w\\}$ such that $u \\in T_{v,i-1}$ and the edge was oriented during the $i$-th round from $w$ to $u$. As $u \\in T_{v,i-1}$, Invariant $(2)$ from the $(i-1)$-th round implies $u \\in S_{i-1}$. This is a contradiction as during the $i$-th round only edges with both endpoints in $U_{i-1}$ can be oriented. Therefore it holds that $T_{v,i} = T_{v,i-1}$, as desired.\n\n \\item Let $v \\in S_i$ be arbitrary. If it also holds that $v \\in S_{i-1}$, then the invariant from round $i-1$ implies that $T_{v,i-1} = T_{v,i}$ contains finitely many vertices and $T_{v,i-1} \\subseteq S_{i-1} \\subseteq S_i$. Thus, it suffices to consider the case that $v \\in S_i \\cap U_{i-1}$. Note that $T_{v,i} \\cap U_i = \\emptyset$. Otherwise there would exist an edge that is oriented away from a vertex in $U_i$ according to $O_i$, but this cannot happen according to the algorithm description. Hence, $T_{v,i} \\subseteq S_i$. Now, for the sake of contradiction, assume that $T_{v,i}$ contains infinitely many vertices. From the definition of $T_{v,i}$, this implies that there exists a sequence of vertices $(v_k)_{k \\geq 1}$ with $v = v_1$ such that for each $k \\geq 1$, $\\{v_{k+1},v_k\\}$ is an edge in $T$ that is oriented from $v_{k+1}$ to $v_k$ according to $O_i$. For each $k \\geq 1$ we have $v_k \\in S_i$. We furthermore know that $v_k \\in U_{i-1}$, as otherwise $T_{v_k,i} = T_{v_k,i-1}$ would contain infinitely many vertices, a contradiction. From the way we orient the edges, a simple induction proof implies that the closest vertex of $v_k$ in the graph $T[U_{i-1}]$ that is contained in $U_i$ has a distance of at least $k$. However, we previously discussed that the distance between $v_k$ and the closest vertex in $U_i$ in the graph $T[U_{i-1}]$ is at most $\\sum_{j=0}^{i} d_j$. This is a contradiction. Hence, $T_{v,i}$ contains finitely many vertices.\n \n \\item It follows directly from the description of how to compute $U_i$ and $H_i$ that the minimum degree of the graph $T[U_i]$ is at least $2$ and that each vertex in $H_i$ has a degree of $\\Delta$ in $T[U_i]$.\n \n \\item Let $u \\in U_i$ be arbitrary. We need to show that $u$ has a distance of at most $\\sum_{j=0}^i d_j$ to the closest vertex in $H_i$ in the graph $T[U_i]$. Let $v$ be the vertex with $u \\in C(v)$. We know that there is no vertex in $H_i$ that is closer to $u$ in $T[U_{i-1}]$ than $v$. Hence, the distance between $u$ and $v$ is at most $\\sum_{j=0}^i d_j$ in the graph $T[U_{i-1}]$. Moreover, from the description of how we compute $U_i$, it follows that all the vertices on the unique path between $u$ and $v$ are contained in $U_i$. Hence, each vertex in $U_i$ has a distance of at most $\\sum_{j=0}^i d_j$ to the closest vertex in $H_i$ in the graph $T[U_i]$, as desired.\n \\end{enumerate}\n\nWe derive an upper bound on the coding radius of the algorithm in three steps. First, for each $i \\in \\mathbb{N}$ and $\\varepsilon > 0$, we derive an upper bound on the coding radius for computing with probability at least $1 - \\varepsilon$ for a given vertex in which of the sets $S_i, U_i$ and $H_i$ it is contained in and for each incident edge whether and how it is oriented in $O_i$, given that each vertex receives as additional input in which of the sets $S_{i-1}, U_{i-1}$ and $H_{i-1}$ it is contained in and for each incident edge whether and how it is oriented in $O_{i-1}$. Given this upper bound on the coding radius, we can use the sequential composition lemma (\\cref{lem:sequential_composition}) and a simple induction proof to give an upper bound on the coding radius for computing with probability at least $1 - \\varepsilon$ for a given vertex in which of the sets $S_i, U_i$ and $H_i$ it is contained in and for each incident edge whether and how it is oriented in $O_i$, this time without providing any additional input. \n\nSecond, we analyze after how many rounds a given vertex is settled with probability at least $1- \\varepsilon$ for a given $\\varepsilon > 0$. \n\nFinally, we combine these two upper bounds to prove \\cref{thm:one_ended_forest}.\n \n \n\\begin{lemma}\n\\label{lem:tail1}\nFor each $i \\in \\mathbb{N}$ and $\\varepsilon \\in (0,0.01]$, let $g_i(\\varepsilon)$ denote the smallest coding radius such that with probability at least $1 - \\varepsilon$ one knows for a given vertex in which of the sets $S_i, U_i$ and $H_i$ it is contained in and for each incident edge whether and how it is oriented in $O_i$, given that each vertex receives as additional input in which of the sets $S_{i-1}, U_{i-1}$ and $H_{i-1}$ it is contained in and for each incident edge whether and how it is oriented in $O_{i-1}$. It holds that $g_i(\\varepsilon) = O(d^2_i \\log (\\Delta\/\\varepsilon))$.\n\\end{lemma}\n\n\\begin{proof}\n When running Ghaffari's uniform MIS algorithm on some graph $G'$ with maximum degree $\\Delta'$, each vertex knows with probability at least $1-\\varepsilon$ whether it is contained in the MIS or not after $O(\\log (\\Delta') + \\log (1\/\\varepsilon))$ communication rounds in $G'$ (\\cite{ghaffari2016MIS}, Theorem 1.1). For the MIS computed during the $i$-th round, $\\Delta' =\\Delta^{O(d_i)}$ and each communication round in $G'$ can be simulated with $O(d_i)$ communication rounds in the tree $T$. Hence, to know for a given vertex with probability at least $1-\\varepsilon$ whether it is contained in the MIS or not, it suffices consider the $O(d_i^2\\log(\\Delta) + d_i\\log(1\/\\varepsilon))$-hop neighborhood around that vertex. Now, let $u$ be an arbitrary vertex. In order to compute in which of the sets $S_i, U_i$ and $H_i$ $u$ is contained in and for each incident edge of $u$ whether and how it is oriented in $O_i$, we not only need to know whether $u$ is in the MIS or not. However, it suffices if we know for all the vertices in the $O(d_i)$-hop neighborhood of $u$ whether they are contained in the MIS or not (on top of knowing for each vertex in the neighborhood in which of the sets $S_{i-1}, U_{i-1}$ and $H_{i-1}$ it is contained in and for each incident edge whether and how it is oriented in $O_{i-1}$). Hence, by a simple union bound over the at most $\\Delta^{O(d_i)}$ vertices in the $O(d_i)$-hop neighborhood around $u$, we obtain $g_i(\\varepsilon) = O(d_i) + O(d_i^2\\log(\\Delta) + d_i\\log(\\Delta^{O(d_i)}\/\\varepsilon)) = O(d_i^2 \\log(\\Delta\/\\varepsilon))$.\n\\end{proof}\n\n\\begin{lemma}\n\\label{lem:tail2}\nFor each $i \\in \\mathbb{N}$ and $\\varepsilon \\in (0,0.01]$, let $h_i(\\varepsilon)$ denote the smallest coding radius such that with probability at least $1 - \\varepsilon$ one knows for a given vertex in which of the sets $S_i, U_i$ and $H_i$ it is contained in and for each incident edge whether and how it is oriented in $O_i$. Then, there exists a constant $c$ independent of $\\Delta$ such that $h_i(\\varepsilon) \\leq 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta\/\\varepsilon)$.\n\\end{lemma}\n\n\\begin{proof}\nBy prove the statement by induction on $i$. For a large enough constant $c$, it holds that $h_1(\\varepsilon) \\leq 2^{2^{3}} \\cdot(c\\log(\\Delta)) \\log(\\Delta\/\\varepsilon)$. Now, consider some arbitrary $i$ and assume that $h_i(\\varepsilon) \\leq 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta\/\\varepsilon)$ for some large enough constant $c$. We show that this implies $h_{i+1}(\\varepsilon) \\leq 2^{2^{(i+1)+2}} \\cdot(c\\log(\\Delta))^{i+1} \\log(\\Delta\/\\varepsilon)$.\nBy the sequential composition lemma (\\cref{lem:sequential_composition}) and assuming that $c$ is large enough, we have\n\n\\begin{align*}\n h_{i+1}(\\varepsilon) &\\leq h_{i}((\\varepsilon\/2)\/\\Delta^{g_{i+1}(\\varepsilon\/2) + 1}) + g_{i+1}(\\varepsilon\/2) \\\\\n &\\leq 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta \\cdot \\Delta^{g_{i+1}(\\varepsilon\/2) + 1}\/(\\varepsilon\/2)) + g_{i+1}(\\varepsilon\/2) \\\\\n &\\leq 2 \\cdot 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta \\cdot \\Delta^{g_{i+1}(\\varepsilon\/2) + 1}\/(\\varepsilon\/2)) \\\\\n &\\leq 2 \\cdot 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta^{(c\/10) \\cdot 2^{2^{i+2}}\\log(\\Delta\/\\varepsilon)}\/(\\varepsilon\/2)) \\\\\n &\\leq 4 \\cdot 2^{2^{i+2}} \\cdot(c\\log(\\Delta))^i \\log(\\Delta^{(c\/10) \\cdot 2^{2^{i+2}}\\log(\\Delta\/\\varepsilon)}) \\\\\n &\\leq (4\/10) \\cdot 2^{2^{i+2}} \\cdot 2^{2^{i+2}} (c \\log(\\Delta))^{i+1} \\log(\\Delta\/\\varepsilon) \\\\\n &\\leq 2^{2^{(i+1)+2}} \\cdot(c\\log(\\Delta))^{i+1} \\log(\\Delta\/\\varepsilon),\n\\end{align*}\n\nas desired.\n\n\\end{proof}\n\n\n \n \n\n\\begin{lemma}\\label{lm:tail0}\nLet $u$ be an arbitrary vertex. For each $\\varepsilon \\in (0,0.01]$, let $f(\\varepsilon)$ denote the smallest $i \\in \\mathbb{N}$ such that $u$ is settled after the $i$-th round with probability at least $1 - \\varepsilon$. There exists a fixed $c \\in \\mathbb{R}$ independent of $\\Delta$ such that $f(\\varepsilon) \\leq \\lceil 1 + \\log \\log \\frac{1}{c} \\log_\\Delta(1\/\\varepsilon)\\rceil$.\n\\end{lemma}\n\n\\begin{proof}\n\nLet $i \\in \\mathbb{N}$ be arbitrary. We show that a given vertex is settled after the $i$-th round with probability at least $1 - O(1\/\\Delta^{\\Omega(d_i\/d_{i-1})})$. For the sake of analysis, we run the algorithm on a finite $\\Delta$-regular high-girth graph instead of an infinite $\\Delta$-regular tree. For now, we additionally assume that no vertex realizes that we don't run the algorithm on an infinite $\\Delta$-regular tree. That is, we assume that for each vertex the coding radius to compute all its local information after the $i$-th round is much smaller than the girth of the graph. \n\nWith this assumption, we give a deterministic upper bound on the fraction of vertices that are not settled after the $i$-th round. On the one hand, the number of vertices that are not settled after the $i$-th round can be upper bounded by $|H_i| \\cdot \\Delta \\cdot O(d_i)$. On the other hand, we will show that the fraction of vertices that are contained in $H_i$ is upper bounded by $1\/\\Delta^{\\Omega(d_i\/d_{i-1})}$. We do this by showing that $C(v)$ contains $\\Delta^{\\Omega(d_i\/d_{i-1})}$ many vertices for a given $v \\in H_i$. Combining these two bounds directly implies that the fraction of unsettled vertices is smaller than \n\n\\[\\frac{\\Delta \\cdot O(d_i)}{\\Delta^{\\Omega(d_i\/d_{i-1})}} = 1\/\\Delta^{\\Omega(d_i\/d_{i-1})}.\\]\n\n\nLet $v \\in H_i$ be arbitrary. We show that $|C(v)| = \\Delta^{\\Omega(d_i\/d_{i-1})}$. Using Invariants (3) and (4) together with a simple induction argument, one can show that there are at least $(\\Delta-1)^{\\lfloor D \/ ((2\\sum_{j=0}^{i-1}d_{j}) + 1)\\rfloor}$ vertices contained in $H_{i-1}$ and whose distance to $v$ is at most $D$ in the graph induced by the vertices in $U_{i-1}$. Furthermore, from the way we defined the clustering $C(v)$ and the fact that two vertices in $H_i$ have a distance of at least $d_i$ in the graph $T[U_{i-1}]$, it follows that all vertices in $U_{i-1}$ having a distance of at most $d_i\/2 -1$ to $v$ are contained in the cluster $C(v)$. Hence, the total number of vertices in $C(v)$ is at least $(\\Delta-1)^{\\lfloor (d_i\/2 - 1) \/ ((2\\sum_{j=0}^{i-1}d_{j}) + 1)\\rfloor} = \\Delta^{\\Omega(d_i\/d_{i-1})}$, as promised.\n\nNow we remove the assumption that no vertex realizes that we don't run the algorithm on an infinite $\\Delta$-regular tree. If a vertex does realize that we don't run the algorithm on an infinite $\\Delta$-regular tree, then we consider the vertex as being unsettled after the $i$-th round. By considering graphs with increasing girth, we can make the expected fraction of vertices that realize that we are not on an infinite $\\Delta$-regular tree arbitrarily small. Combining this observation with the previous discussion, this implies that the expected fraction of vertices that are not settled after the $i$-th round is at most $1\/\\Delta^{\\Omega(d_i\/d_{i-1})}$. By symmetry, each vertex has the same probability of being settled after the $i$-th round. Hence, the probability that a given vertex is settled after the $i$-th round when run on a graph with sufficiently large girth is $1 - 1\/\\Delta^{\\Omega(d_i\/d_{i-1})}$, and the same holds for each vertex when we run the algorithm on an infinite $\\Delta$-regular tree. Thus, there exists a constant $c$ such that the probability that a given vertex is unsettled after the $i$-th round is at most $1\/\\Delta^{c \\cdot d_i\/d_{i-1}} = 1\/\\Delta^{c \\cdot 2^{2^{i-1}}}$. Setting $i = \\lceil 1 + \\log \\log \\frac{1}{c} \\log_\\Delta(1\/\\varepsilon)\\rceil$ finishes the proof.\n\\end{proof}\n\n\n\n\n\nWe are now finally ready to finish the proof of \\cref{thm:one_ended_forest}. Let $u$ be an arbitrary vertex and $\\varepsilon \\in (0,0.01]$. We need to compute an upper bound on the coding radius that is necessary for $u$ to know all the vertices in $\\mathcal{T}^{\\leftarrow}(u)$ with probability at least $1-\\varepsilon$. To compute such an upper bound for the required coding radius it suffices to find an $i^*$ and an $R^*$ such that the following holds. First, $u$ is settled after $i^*$ rounds with probability at least $1 - \\varepsilon\/2$. Second, $u$ knows for each edge in its $O(d_{i^*})$-hop neighborhood whether it is oriented in the partial orientation $O_{i^*}$, and if yes, in which direction, by only considering its $R^*$-hop neighborhood with probability at least $1 - \\varepsilon\/2$. By a union bound, both of these events occur with probability at least $1-\\varepsilon$. Moreover, if both events occur $u$ knows all the vertices in the set $\\mathcal{T}^{\\leftarrow}(u)$. The reason is as follows. If $u$ is settled after round $i^*$, then it follows from the previous analysis that only vertices in the $O(d_{i^*})$-hop neighborhood of $u$ can be contained in $\\mathcal{T}^{\\leftarrow}(u)$. Moreover, for each vertex in its $O(d_{i^*})$-hop neighborhood, $u$ can determine if the vertex is contained in $\\mathcal{T}^{\\leftarrow}(u)$ if it knows all the edge orientations on the unique path between itself and that vertex after the $i^*$-th round. Hence, it remains to find concrete values for $i^*$ and $R^*$. According to \\cref{lm:tail0}, we can choose $i^* = \\lceil 1 + \\log \\log \\frac{1}{c} \\log_\\Delta (2\/\\varepsilon)\\rceil$ for some large enough constant $c$. Moreover, it follows from a union bound that all vertices in the $O(d_{i^*})$-hop neighborhood around $u$ know with probability at least $1-\\varepsilon\/2$ the orientation of all its incident edges according to $O_{i^*}$ by only considering their $h_{i^*}((\\varepsilon\/2)\/\\Delta^{O(d_{i^*})})$-hop neighborhood. Hence, we can set \n\n\\begin{align*}\nR^* &= O(d_{i^*}) + h_{i^*}((\\varepsilon\/2)\/\\Delta^{O(d_{i^*})}) \\\\\n &\\leq O(d_{i^*}) + 2^{2^{i^* + 2}} \\cdot (c \\log(\\Delta))^{i^*}\\log(\\Delta^{O(d_{i^*})}\/\\varepsilon) \\\\\n &= O(2^{2^{i^*}}\\cdot 2^{2^{i^* + 2}} (c \\log(\\Delta))^{i^*}\\log(\\Delta\/\\varepsilon)) \\\\\n &= O(2^{2^{i^* + 3}} (c \\log(\\Delta))^{i^*}\\log(\\Delta\/\\varepsilon)) \\\\\n &= O(2^{2^{\\log \\log \\frac{1}{c} \\log (2\/\\varepsilon) + 5}} (c \\log(\\Delta))^{\\log \\log \\frac{1}{c} \\log (2\/\\varepsilon) + 2}\\log(\\Delta\/\\varepsilon)) \\\\\n &= \\log(1\/\\varepsilon)^{32} \\cdot \\log(\\Delta\/\\varepsilon) \\cdot \\log\\log(1\/\\varepsilon)^{O(\\log\\log(\\Delta))}.\n\\end{align*}\n\nAs $\\Delta = O(1)$, it therefore holds that\n\n\\[P(R(v)) = \\operatorname{poly}(\\log(1\/\\varepsilon)) \\geq 1 - \\varepsilon, \\]\n\nas desired.\n\n\n\n\n\n\n\\section{$\\mathsf{BAIRE} = \\mathsf{LOCAL}(O(\\log n))$}\n\\label{sec:baire}\n\nIn this section, we show that on $\\Delta$-regular trees the classes $\\mathsf{BAIRE}$ and $\\mathsf{LOCAL}(O(\\log(n)))$ are the same.\nAt first glance, this result looks rather counter-intuitive.\nThis is because in finite $\\Delta$-regular trees every vertex can see a leaf of distance $O(\\log(n))$, while there are no leaves at all in an infinite $\\Delta$-regular tree.\nHowever, there is an intuitive reasons why these classes are the same: in both setups there is a technique to decompose an input graph into a hierarchy of subsets. Furthermore, the existence of a solution that is defined inductively with respect to these decompositions can be characterized by the same combinatorial condition of Bernshteyn \\cite{Bernshteyn_work_in_progress}.\nWe start with a high-level overview of the decomposition techniques used in both contexts.\n\n\n\\paragraph{{\\sf Rake}\\ and {\\sf Compress}}\nThe hierarchical decomposition in the context of distributed computing is based on a variant of a decomposition algorithm of\nMiller and Reif~\\cite{MillerR89}.\nTheir original decomposition algorithm works as follows.\nStart with a tree $T$, and repeatedly apply the following two operations alternately: {\\sf Rake}\\ (remove all degree-1 vertices) and {\\sf Compress}\\ (remove all degree-2 vertices).\nThen $O(\\log n)$ iterations suffice to remove all vertices in $T$~\\cite{MillerR89}. To view it another way, this produces a decomposition of the vertex set $V$ into $2L - 1$ layers \\[V = \\VR{1} \\cup \\VC{1} \\cup \\VR{2} \\cup \\VC{2} \\cup \\VR{3} \\cup \\VC{3} \\cup \\cdots \\cup \\VR{L},\\] \nwith $L= O(\\log n)$, where $\\VR{i}$ is the set of vertices removed during the $i$-th {\\sf Rake}\\ operation and $\\VC{i}$ is the set of vertices removed during the $i$-th {\\sf Compress}\\ operation. We will use a variant~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} of this decomposition in the proof of \\cref{thm:ellfull_to_logn}.\n\n\n\n\nVariants of this decomposition turned out to be useful in designing $\\mathsf{LOCAL}$ algorithms~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup,ChangHLPU20,chang2020n1k_speedups}. \nIn our context, we assume that the given LCL satisfies a certain combinatorial condition and then find a solution inductively, in the reversed order of the construction of the decomposition.\nNamely, in the {\\sf Rake}\\ step we want to be able to existentially extend the inductive partial solution to all relative degree $1$-vertices (each $v \\in \\VR{i}$ has degree at most 1 in the subgraph induced by $\\VR{i} \\cup \\cdots \\cup \\VR{L}$) and in the {\\sf Compress}\\ step we want to extend the inductive partial solution to paths with endpoints labeled from the induction (the vertices in $\\VC{i}$ form degree-2 paths in the subgraph induced by $\\VC{i} \\cup \\cdots \\cup \\VR{L}$).\n\n\n\\paragraph{$\\TOAST$}\nFinding a hierarchical decomposition in the context of descriptive combinatorics is tightly connected with the notion of \\emph{Borel hyperfiniteness}.\nUnderstanding what Borel graphs are Borel hyperfinite is a major theme in descriptive set theory \\cite{doughertyjacksonkechris,gaojackson,conley2020borel}.\nIt is known that grids, and generally polynomial growth graphs are hyperfinite, while, e.g., acyclic graphs are not in general hyperfinite \\cite{jackson2002countable}.\nA strengthening of hyperfiniteness that is of interest to us is called a \\emph{toast} \\cite{gao2015forcing,conleymillerbound}.\nA $q$-toast, where $q\\in \\mathbb{N}$, of a graph $G$ is a collection $\\mathcal{D}$ of fintie subsets of $G$ with the property that (i) every pair of vertices is covered by an element of $\\mathcal{D}$ and (ii) the boundaries of every $D\\not=E\\in \\mathcal{D}$ are at least $q$ apart.\nThe idea to use a toast structure to solve LCLs appears in \\cite{conleymillerbound} and has many applications since then \\cite{gao2015forcing,marksunger}.\nThis approach has been formalized in \\cite{grebik_rozhon2021toasts_and_tails}, where the authors introduce $\\TOAST$ algorithms.\nRoughly speaking, an LCL $\\Pi$ admits a $\\TOAST$ algorithm if there is $q\\in \\mathbb{N}$ and a partial extending function (the function is given a finite subset of a tree that is partially colored and outputs an extension of this coloring on the whole finite subset) that has the property that whenever it is applied inductively to a $q$-toast, then it produces a $\\Pi$-coloring.\nAn advantage of this approach is that once we know that a given Borel graph admits, e.g., a Borel toast structure and a given LCL $\\Pi$ admits a $\\TOAST$ algorithm, then we may conclude that $\\Pi$ is in the class $\\mathsf{BOREL}$.\nSimilarly for $\\mathsf{MEASURE}$, $\\mathsf{BAIRE}$ or $\\mathsf{OLOCAL}$, we refer the reader to \\cite{grebik_rozhon2021toasts_and_tails} for more details and results concerning grids.\n\nIn the case of trees there is no way of constructing a Borel toast in general, however, it is a result of Hjorth and Kechris \\cite{hjorth1996borel} that every Borel graph is hyperfinite on a comeager set for every compatible Polish topology.\nA direct consequence of \\cite[Lemma~3.1]{marks2016baire} together with a standard construction of toast via Voronoi cells gives the following strengthening to toast.\nWe include a sketch of the proof for completeness.\n\n\\begin{restatable}{proposition}{BaireToast}\n\\label{pr:BaireToast}\nLet $\\mathcal{G}$ be a Borel graph on a Polish space $(X,\\tau)$ with degree bounded by $\\Delta\\in \\mathbb{N}$.\nThen for every $q>0$ there is a Borel $\\mathcal{G}$-invariant $\\tau$-comeager set $C$ on which $\\mathcal{G}$ admits a Borel $q$-toast.\n\\end{restatable}\n\\begin{proof}[Proof sketch]\nLet $\\{A_n\\}_{n\\in \\mathbb{N}}$ be a sequence of MIS with parameter $f(n)$ as in \\cite[Lemma~3.1]{marks2016baire} for a sufficiently fast growing function $f(n)$, e.g., $f(n)=(2q)^{n^2}$.\nThen $C=\\bigcup_{n\\in \\mathbb{N}}A_n$ is a Borel $\\tau$-comeager set that is $\\mathcal{G}$-invariant.\nWe produce a toast structure in a standard way, e.g., see~\\cite[Appendix~A]{grebik_rozhon2021toasts_and_tails}.\n\nLet $B_n(x)$ denote the ball of radius $f(n)\/3$ and $R_n(x)$ the ball of radius $f(n)\/4$ around $x\\in A_n$.\nIteratively, define cells $\\mathcal{D}_n=\\{C_{n}(x)\\}_{x\\in A_n}$ as follows.\nSet $C_1(x)=R_1(x)$.\nSuppose that $\\mathcal{D}_n$ has been defined and set\n\\begin{itemize}\n\t\\item $H^{n+1}(x,n+1):=R_{n+1}(x)$ for every $x\\in A_{n+1}$,\n\t\\item if $1\\le i\\le n$ and $\\left\\{H^{n+1}(x,i+1)\\right\\}_{x\\in A_{k+1}}$ has been defined, then we put\n\t$$H^{n+1}(x,i)=\\bigcup \\left\\{B_i(y):H^{n+1}(x,i+1)\\cap B_i(y)\\not=\\emptyset\\right\\}$$\n\tfor every $x\\in A_{n+1}$,\n\t\\item set $C_{n+1}(x):=H^{n+1}(x,1)$ for every $x\\in A_{n+1}$, this defines $\\mathcal{D}_{n+1}$.\n\\end{itemize}\nThe fact that $\\mathcal{D}=\\bigcup_{n\\in\\mathbb{N}} \\mathcal{D}_n$ is a $q$-toast on $C$ follows from the fact that $C=\\bigcup_{n\\in \\mathbb{N}}A_n$ together with the fact that the boundaries are $q$ separated which can be shown as in \\cite[Appendix~A]{grebik_rozhon2021toasts_and_tails}.\n\\end{proof}\n\n\n\nTherefore to understand LCLs in the class $\\mathsf{BAIRE}$ we need understand what LCLs on trees admit $\\TOAST$ algorithm.\nIt turns out that these notions are equivalent, again by using the combinatorial characterization of Bernshteyn \\cite{Bernshteyn_work_in_progress} that we now discuss.\n\n\n\\paragraph{Combinatorial Condition -- $\\ell$-full set}\nIn both decompositions, described above, we need to extend a partial coloring along paths that have their endpoints colored from the inductive step.\nThe precise formulation of the combinatorial condition that captures this demand was extracted by Bernshteyn \\cite{Bernshteyn_work_in_progress}.\nHe proved that it characterizes the class $\\mathsf{BAIRE}$ for Cayley graphs of virtually free groups.\nNote that this class contains, e.g., $\\Delta$-regular trees with a proper edge $\\Delta$-coloring.\n\n\\begin{definition}[Combinatorial condition -- an $\\ell$-full set]\n\\label{def:ellfull}\nLet $\\Pi = (\\Sigma, \\mathcal{V}, \\mathcal{E})$ be an LCL and $\\ell\\geq 2$.\nA set $\\mathcal{V}' \\subseteq \\mathcal{V}$ is \\emph{$\\ell$-full} whenever the following is satisfied.\nTake a path with at least $\\ell$ vertices, and add half-edges to it so that each vertex has degree $\\Delta$.\nTake any $c_1, c_2 \\in \\mathcal{V}'$ and label arbitrarily the half-edges around the endpoints with $c_1$ and $c_2$, respectively.\nThen there is a way to label the half-edges around the remaining $\\ell-2$ vertices with configurations from $\\mathcal{V}'$ such that all the $\\ell-1$ edges on the path have valid edge configuration on them. \n\\end{definition}\n\nNow we are ready to formulate the result that combines Bernshteyn's result \\cite{Bernshteyn_work_in_progress} (equivalence between (1.) and (2.), and the moreover part) with the main results of this section.\nThis also shows the remaining implications in \\cref{fig:big_picture_trees}.\n\n\\begin{theorem}\\label{thm:MainBaireLog}\nLet $\\Pi$ be an LCL on regular trees. \nThen the following are equivalent:\n\\begin{enumerate}\n \\item $\\Pi\\in \\mathsf{BAIRE}$,\n \\item $\\Pi$ admits a $\\TOAST$ algorithm,\n \\item $\\Pi$ admits an $\\ell$-full set,\n \\item $\\Pi\\in \\mathsf{LOCAL}(O(\\log(n)))$.\n\\end{enumerate}\nMoreover, any of the equivalent conditions is necessary for $\\Pi\\in \\mathsf{fiid}$.\n\\end{theorem}\n\nNext we discuss the proof of \\cref{thm:MainBaireLog}.\nWe refer the reader to Bernshteyn's paper \\cite{Bernshteyn_work_in_progress} for full proofs in the case of $\\mathsf{BAIRE}$ and $\\mathsf{fiid}$, here we only sketch the argument for completeness.\nWe also note that instead of using the toast construction, he used a path decomposition of acyclic graphs of Conley, Marks and Unger \\cite{conley2020measurable}.\n\n\n\\subsection{Sufficiency}\n\nWe start by showing that the combinatorial condition is sufficient for $\\mathsf{BAIRE}$ and $\\mathsf{LOCAL}(O(\\log(n)))$.\nNamely, it follows from the next results together with \\cref{pr:BaireToast} that (2.) implies all the other conditions in \\cref{thm:MainBaireLog}.\nAs discussed above the main idea is to color inductively along the decompositions.\n\n\\begin{proposition}\\label{pr:toastable}\nLet $\\Pi=(\\Sigma, \\mathcal{V}, \\mathcal{E})$ be an LCL that admits $\\ell$-full set $\\mathcal{V}'\\subseteq \\mathcal{V}$ for some $\\ell>0$.\nThen $\\Pi$ admits a $\\TOAST$ algorithm that produces a $\\Pi$-coloring for every $(2\\ell+2)$-toast $\\mathcal{D}$.\n\\end{proposition}\n\\begin{proof}[Proof sketch]\nOur aim is to build a partial extending function.\nSet $q:=2\\ell+2$.\nLet $E$ be a piece in a $q$-toast $\\mathcal{D}$ and suppose that $D_1,\\dots,D_k\\in \\mathcal{D}$ are subsets of $E$ such that the boundaries are separated.\nSuppose, moreover, that we have defined inductively a coloring of half-edges of vertices in $D=\\bigcup D_i$ using only vertex configurations from $\\mathcal{V}'$ such that every edge configuration $\\mathcal{E}$ is satisfied for every edge in $D$.\n\nWe handle each connected component of $E\\setminus D$ separately.\nLet $A$ be one of them.\nLet $u\\in A$ be a boundary vertex of $E$.\nSuch an vertex exists since every vertex in $E$ has degree $\\Delta$.\nThe distance of $u$ and any $D_i$ is at least $2\\ell+2$ for every $i\\in [k]$.\nWe orient all the edges from $A$ towards $u$.\nMoreover if $v_i\\in A$ is a boundary vertex of some $D_i$ we assign to $v_i$ a path $V_i$ of length $\\ell$ towards $u$.\nNote that $V_i$ and $V_j$ have distance at least $1$, in particular, are disjoint for $i\\not=j\\in [k]$ .\nNow, until you encounter some path $V_i$, color any in manner half-edges of vertices in $A$ inductively starting at $u$ in such a way that edge configurations $\\mathcal{E}$ are satisfied on every edge and only vertex configurations from $\\mathcal{V}'$ are used.\nUse the definition of $\\ell$-full set to find a coloring of any such $V_i$ and continue in a similar manner until the whole $A$ is colored.\n\\end{proof}\n\n\n\n\n\\begin{proposition}[$\\ell$-full $\\Rightarrow$ $\\mathsf{LOCAL}(O(\\log n))$]\\label{thm:ellfull_to_logn}\nLet $\\Pi = (\\Sigma, \\mathcal{V}, \\mathcal{E})$ be an LCL with an $\\ell$-full set $\\mathcal{V}' \\subseteq \\mathcal{V}$. Then $\\Pi$ can be solved in $O(\\log n)$ rounds in $\\mathsf{LOCAL}$.\n\\end{proposition}\n\\begin{proof}\nThe proof uses a variant of the rake-and-compress decomposition considered in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}.\n\n\\paragraph{The Decomposition}\nThe decomposition is parameterized an integer $\\ell' \\geq 1$, and it decomposes the vertices of $T$ into $2L - 1$ layers \\[V = \\VR{1} \\cup \\VC{1} \\cup \\VR{2} \\cup \\VC{2} \\cup \\VR{3} \\cup \\VC{3} \\cup \\cdots \\cup \\VR{L},\\] \nwith $L= O(\\log n)$. We write $\\GC{i}$ to denote the subtree induced by the vertices $\\left(\\bigcup_{j=i+1}^{L} \\VR{j}\\right) \\cup \\left( \\bigcup_{j=i}^{L-1} \\VC{j}\\right)$. Similarly, $\\GR{i}$ is the subtree induced by the vertices \n$\\left(\\bigcup_{j=i}^L \\VR{j}\\right) \\cup \\left( \\bigcup_{j=i}^{L-1} \\VC{j}\\right)$.\nThe sets $\\VR{i}$ and $\\VC{i}$ are required to satisfy the following requirements. \n\\begin{itemize}\n \\item Each $v \\in \\VR{i}$ has degree at most one in the graph $\\GR{i}$.\n \\item Each $v \\in \\VC{i}$ has degree exactly two in the graph $\\GC{i}$. Moreover, the $\\VC{i}$-vertices in $\\GC{i}$ form paths with $s$ vertices, with $\\ell' \\leq s \\leq 2 \\ell'$.\n\\end{itemize}\nFor any given constant $\\ell' \\geq 1$, it was shown in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} that such a decomposition of a tree $T$ can be computed in $O(\\log n)$ rounds. See \\cref{fig:rake_and_compress} for an example of such a decomposition with $\\ell' = 4$.\n\n\\paragraph{The Algorithm}\nGiven such a decomposition with $\\ell' = \\max\\{1, \\ell - 2\\}$, $\\Pi$ can be solved in $O(\\log n)$ rounds by labeling the vertices in this order: $\\VR{L}$, $\\VC{L-1}$, $\\VR{L-1}$, $\\ldots$, $\\VR{1}$, as follows. The algorithm only uses the vertex configurations in the $\\ell$-full set $\\mathcal{V}'$.\n\n\\paragraph{Labeling $\\VR{i}$}\nSuppose all vertices in $\\VR{L}, \\VC{L-1}, \\VR{L-1}, \\ldots, \\VC{i}$ have been labeled using $\\mathcal{V}'$. Recall that each $v \\in \\VR{i}$ has degree at most one in the graph $\\GR{i}$.\nIf $v \\in \\VR{i}$ has no neighbor in $\\VR{L} \\cup \\VC{L-1} \\cup \\VR{L-1} \\cup \\cdots \\cup \\VC{i}$, then we can label the half edges surrounding $v$ by any $c \\in \\mathcal{V}'$.\nOtherwise, $v \\in \\VR{i}$ has exactly one neighbor $u$ in $\\VR{L} \\cup \\VC{L-1} \\cup \\VR{L-1} \\cup \\cdots \\cup \\VC{i}$. Suppose the vertex configuration of $u$ is $c$, where the half-edge label on $\\{u,v\\}$ is $\\mathtt{a} \\in c$. A simple observation from the definition of $\\ell$-full sets is that for any $c \\in \\mathcal{V}'$ and any $\\mathtt{a} \\in c$, there exist $c' \\in \\mathcal{V}'$ and $\\mathtt{a}' \\in c'$ in such a way that $\\{\\mathtt{a}, \\mathtt{a}'\\} \\in \\mathcal{E}$.\nHence we can label the half edges surrounding $v$ by $c' \\in \\mathcal{V}'$ where the half-edge label on $\\{u,v\\}$ is $\\mathtt{a}' \\in c'$. \n\n\n\\paragraph{Labeling $\\VC{i}$} Suppose all vertices in $\\VR{L}, \\VC{L-1}, \\VR{L-1}, \\ldots, \\VR{i+1}$ have been labeled using $\\mathcal{V}'$. Recall that the $\\VC{i}$-vertices in $\\GC{i}$ form degree-2 paths $P=(v_1, v_2, \\ldots, v_s)$, with $\\ell' \\leq s \\leq 2 \\ell'$.\nLet $P'= (x, v_1, v_2, \\ldots, v_s, y)$ be the path resulting from appending to $P$ the neighbors of the two end-points of $P$ in $\\GC{i}$. The two vertices $x$ and $y$ are in $\\VR{L} \\cup \\VC{L-1} \\cup \\VR{L-1} \\cup \\cdots \\cup \\VR{i+1}$, so they have been assigned half-edge labels using $\\mathcal{V}'$. Since $P'$ contains at least $\\ell'+2 \\geq \\ell$ vertices, the definition of $\\ell$-full sets ensures that we can label $v_1, v_2, \\ldots, v_s$ using vertex configurations in $\\mathcal{V}'$ in such a way that the half-edge labels on $\\{x, v_1\\}, \\{v_1, v_2\\}, \\ldots, \\{v_s, y\\}$ are all in $\\mathcal{E}$.\n\\end{proof}\n\n\n\n\n\n\\begin{figure}[h!]\n \\centering\n \\includegraphics[width= .8\\textwidth]{fig\/y1.pdf}\n \\caption{The variant of the rake-and-compress decomposition used in the proof of \\cref{thm:ellfull_to_logn}. }\n \\label{fig:rake_and_compress}\n\\end{figure}\n\n\n\n\\subsection{Necessity}\n\nWe start by sketching that (2.) in \\cref{thm:MainBaireLog} is necessary for $\\mathsf{BAIRE}$ and $\\mathsf{fiid}$.\n\n\n\\begin{theorem}[Bernshteyn \\cite{Bernshteyn_work_in_progress}]\nLet $\\Pi=(\\Sigma,\\mathcal{V},\\mathcal{E})$ be an LCL and suppose that $\\Pi\\in \\mathsf{BAIRE}$ or $\\Pi\\in \\mathsf{fiid}$.\nThen $\\Pi$ admits an $\\ell$-full set $\\mathcal{V}'\\subseteq\\mathcal{V}$ for some $\\ell>0$.\n\\end{theorem}\n\\begin{proof}[Proof Sketch]\n\n\nWe start with $\\mathsf{BAIRE}$.\nSuppose that every Borel acyclic $\\Delta$-regular graph admits a Borel solution on a $\\tau$-comeager set for every compatible Polish topology $\\tau$.\nIn particular, this holds for the Borel graph induced by the standard generators of the free product of $\\Delta$-copies of $\\mathbb{Z}_2$ on the free part of the shift action on the alphabet $\\{0,1\\}$ endowed with the product topology.\nLet $F$ be such a solution.\nWrite $\\mathcal{V}'\\subseteq \\mathcal{V}$ for the configurations of half-edge labels around vertices that $F$ outputs on a non-meager set.\nLet $C$ be a comeager set on which $F$ is continuous. Then, every element of $\\mathcal{V}'$ is encoded by some finite window in the shift on $C$, that is, for each element there are a $k \\in \\mathbb{N}$ and function $s:B(1,k) \\to \\{0,1\\}$ such that $F$ is constant on the set $N_s \\cap C$ (where $N_s$ is the basic open neighbourhood determined by $s$, and $B(1,k)$ is the $k$-neighbourhood of the identity in the Cayley graph of the group).\nSince $\\mathcal{V}'$ is finite, we can take $t>0$ to be the maximum of such $k$'s.\nIt follows by standard arguments that $\\mathcal{V}'$ is $\\ell$-full for $\\ell>2t+1$.\n\nA similar argument works for the $\\mathsf{fiid}$, however, for the sake of brevity, we sketch a shorter argument that uses the fact that there must be a correlation decay for factors of iid's. Let $\\Pi\\in \\mathsf{fiid}$.\nThat is, there is an $\\operatorname{Aut}(T)$-equivariant measurable function from iid's on $T$ (without colored edges this time) into the space of $\\Pi$-colorings.\nLet $\\mathcal{V}'$ be the set of half-edges configurations around vertices that have non-zero probability to appear.\nLet $u,v\\in T$ be vertices of distance $k_0\\in \\mathbb{N}$.\nBy \\cite{BGHV} the correlation between the configurations around $u$ and $v$ tends to $0$ as $k_0\\to \\infty$.\nThis means that if the distance is big enough, then all possible pairs of $\\mathcal{V}'$ configurations need to appear.\n\\end{proof}\n\nTo finish the proof of \\cref{thm:MainBaireLog} we need to demonstrate the following theorem. Note that $\\mathsf{LOCAL}(n^{o(1)}) = \\mathsf{LOCAL}(O(\\log n))$ according to the $\\omega(\\log n)$ -- $n^{o(1)}$ complexity gap~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}.\n\n\\begin{restatable}{theorem}{logcomb}\\label{thm:logn_to_ellfull}\nLet $\\Pi = (\\Sigma, \\mathcal{V}, \\mathcal{E})$ be an LCL solvable in $\\mathsf{LOCAL}(n^{o(1)})$ rounds.\nThen there exists an $\\ell$-full set $\\mathcal{V}' \\subseteq \\mathcal{V}$ for some $\\ell \\geq 2$.\n\\end{restatable}\n\n\nThe rest of the section is devoted to the proof of \\cref{thm:logn_to_ellfull}. We start with the high-level idea of the proof.\nA natural attempt for showing $\\mathsf{LOCAL}(n^{o(1)})$ $\\Rightarrow$ $\\ell$-full is to simply take any $\\mathsf{LOCAL}(n^{o(1)})$ algorithm $\\mathcal{A}$ solving $\\Pi$, and then take $\\mathcal{V}'$ to be all vertex configurations that can possibly occur in an output of $\\mathcal{A}$. It is not hard to see that this approach does not work in general, because the algorithm might use a special strategy to label vertices with degree smaller than $\\Delta$. Specifically, there might be some vertex configuration $c$ used by $\\mathcal{A}$ so that some $\\mathtt{a} \\in c$ will only be used to label \\emph{virtual} half edges. It will be problematic to include $c$ in $\\mathcal{V}'$.\n\nTo cope with this issue, we do not deal with general bounded-degree trees. Instead, we construct recursively a sequence $(W^\\ast_1, W^\\ast_2, \\ldots, W^\\ast_L)$ of sets of rooted, layered, and \\emph{partially labeled} tree in a special manner. A tree $T$ is included in $W^\\ast_i$ if it can be constructed by gluing a multiset of rooted trees in $W^\\ast_{i-1}$ and a new root vertex $r$ in a certain fixed manner. A vertex is said to be in layer $i$ if it is introduced during the $i$-th step of the construction, i.e., it is introduced as the root $r$ during the construction of $W^\\ast_i$ from $W^\\ast_{i-1}$. All remaining vertices are said to be in layer 0.\n\n\nWe show that each $T \\in W^\\ast_L$ admits a correct labeling that extends the given partial labeling, as these partial labelings are computed by a simulation of $\\mathcal{A}$.\nMoreover, in these correct labelings, the variety of possible configurations of half-edge labels around vertices in different non-zero layers is the same for each layer.\nThis includes vertices of non-zero layer whose half-edges are labeled by the given partial labeling.\nWe simply pick $\\mathcal{V}'$ to be the set of all configurations of half-edge labels around vertices that can appear in a non-zero layer in a correct labeling of a tree $T \\in W^\\ast_L$.\nOur construction ensures that each $c \\in \\mathcal{V}'$ appears as the labeling of some degree-$\\Delta$ vertex in some tree that we consider.\n\n\nThe proof that $\\mathcal{V}'$ is an $\\ell$-full set is based on finding paths using vertices of non-zero layers connecting two vertices with any two vertex configurations in $\\mathcal{V}'$ in different lengths. These paths exist because the way rooted trees in $W^\\ast_{i-1}$ are glued together in the construction of $W^\\ast_{i}$ is sufficiently flexible.\nThe reason that we need $\\mathcal{A}$ to have complexity $\\mathsf{LOCAL}(n^{o(1)})$ is that the construction of the trees can be parameterized by a number $w$ so that all the trees have size polynomial in $w$ and the vertices needed to be assigned labeling are at least distance $w$ apart from each other. Since the number of rounds of $\\mathcal{A}$ executed on trees of size $w^{O(1)}$ is much less than $w$, each labeling assignment can be calculated locally and independently. \nThe construction of the trees as well as the analysis are based on a machinery developed in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}.\n Specifically, we will consider the equivalence relation $\\overset{\\star}{\\sim}$ defined in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} and prove some of its properties, including a pumping lemma for bipolar trees. The exact definition of $\\overset{\\star}{\\sim}$ in this paper is different from the one in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup} because the mathematical formalism describing LCL problems in this paper is different from the one in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}.\nAfter that, we will consider a procedure for gluing trees parameterized by a labeling function $f$ similar to the one used in~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup}. \nWe will apply this procedure iteratively to generate a set of trees. We will show that the desired $\\ell$-full set $\\mathcal{V}' \\subseteq \\mathcal{V}$ can be constructed by considering the set of all possible correct labeling of these trees. \n \n\n\\paragraph{The Equivalence Relation $\\overset{\\star}{\\sim}$} \nWe consider trees with a list of designated vertices $v_1, v_2, \\ldots, v_k$ called \\emph{poles}. A \\emph{rooted tree} is a tree with one pole $r$, and a \\emph{bipolar tree} is a tree with two poles $s$ and $t$. \nFor a tree $T$ with its poles $S=(v_1, v_2, \\ldots, v_k)$ with $\\deg(v_i) = d_i < \\Delta$, we denote by $h_{(T,S)}$ the function that maps each choice of the \\emph{virtual} half-edge labeling surrounding the poles of $T$ to $\\textsf{YES}$ or $\\textsf{NO}$, indicating whether such a partial labeling can be completed into a correct complete labeling of $T$. \nMore specifically, consider \n\\[X=(\\mathcal{I}_1, \\mathcal{I}_2, \\ldots, \\mathcal{I}_k),\\]\nwhere $\\mathcal{I}_i$ is a size-$(\\Delta - d_i)$ multiset of labels in $\\Sigma$, for each $1 \\leq i \\leq k$. Then $h_{(T,S)}(X) = \\textsf{YES}$ if there is a correct labeling of $T$ such that the $\\Delta - d_i$ virtual half-edge labels surrounding $v_i$ are labeled by $\\mathcal{I}_i$, for each $1 \\leq i \\leq k$. This definition can be generalized to the case $T$ is already partially labeled in the sense that some of the half-edge labels have been fixed. In this case, $h_{(T,S)}(X) = \\textsf{NO}$ whenever $X$ is incompatible with the given partial labeling.\n\nLet $T$ be a tree with poles $S = (v_1, v_2, \\ldots, v_k)$ and let $T'$ be another tree with poles $S' = (v_1', v_2', \\ldots, v_k')$ such that $\\deg(v_i) = \\deg(v_i') = d_i$ for each $1 \\leq i \\leq k$.\nThen we write $T_1 \\overset{\\star}{\\sim} T_2$ if $h_{(T,S)} = h_{(T',S')}$.\n\nGiven an LCL problem $\\Pi$, it is clear that the number of equivalence classes of rooted trees and bipolar trees w.r.t.~$\\overset{\\star}{\\sim}$ is finite.\nFor a rooted tree $T$, denote by $\\textsf{Class}_1(T)$ the equivalence class of $T$. For a bipolar tree $H$, denote by $\\texttt{type}(H)$ the equivalence class of $H$.\n\n\n\n\n\\paragraph{Subtree Replacement}\nThe following lemma provides a sufficient condition that the equivalence class of a tree $T$ is invariant of the equivalence class of its subtree $T'$.\nWe note that a real half edge in $T$ might become virtual in its subtree $T'$.\nConsider a vertex $v$ in $T'$ and its neighbor $u$ that is in $T$ but not in $T'$. Then the half edge $(v, \\{u,v\\})$ is real in $T$ and virtual in $T'$.\n\n\n\\begin{lemma}[Replacing subtrees]\\label{lem:replace}\nLet $T$ be a tree with poles $S$.\nLet $T'$ be a connected subtree of $T$ induced by $U \\subseteq V$, where $V$ is the set of vertices in $T$.\nWe identify a list of designated vertices $S' = (v_1', v_2', \\ldots, v_k')$ in $U$ satisfying the following two conditions to be the poles of $T'$.\n\\begin{itemize}\n \\item $S \\cap U \\subseteq S'$.\n \\item Each edge $e=\\{u,v\\}$ connecting $u \\in U$ and $ v \\in V \\setminus U$ must satisfy $u \\in S'$.\n\\end{itemize}\nLet $T''$ be another tree with poles $S'' = (v_1'', v_2'', \\ldots, v_k'')$ that is in the same equivalence class as $T'$.\nLet $T^\\ast$ be the result of replacing $T'$ by $T''$ in $T$ by identifying $v_i' = v_i''$ for each $1 \\leq i \\leq k$.\nThen $T^\\ast$ is in the same equivalence class as $T$. \n\\end{lemma}\n\\begin{proof}\nTo prove the lemma, by symmetry, it suffices to show that starting from any correct labeling $\\mathcal{L}$ of $T$, it is possible to find a correct labeling $\\mathcal{L}^\\ast$ of $T^\\ast$ in such a way that the multiset of the\nvirtual half-edge labels surrounding each pole in $S$ remain the same.\n\nSuch a correct labeling $\\mathcal{L}^\\ast$ of $T^\\ast$ is constructed as follows. If $v \\in V \\setminus U$, then we simply adopt the given labeling $\\mathcal{L}$ of $v$ in $T$.\nNext, consider the vertices in $U$.\nSet $X=(\\mathcal{I}_1, \\mathcal{I}_2, \\ldots, \\mathcal{I}_k)$ to be the one compatible with the labeling $\\mathcal{L}$ of $T$ restricted to the subtree $T'$ in the sense that $\\mathcal{I}_i$ is the multiset of the virtual half-edge labels surrounding the pole $v_i'$ of $T'$, for each $1 \\leq i \\leq k$. We must have $h_{T',S'}(X) = \\textsf{YES}$.\nSince $T'$ and $T''$ are in the same equivalence class, we have $h_{T'',S''}(X) = \\textsf{YES}$, and so we can find a correct labeling $\\mathcal{L}''$ of $T''$ that is also compatible with $X$. Combining this labeling $\\mathcal{L}''$ of the vertices $U$ in $T''$ with the labeling $\\mathcal{L}$ of the vertices $V \\setminus U$ in $T$, we obtain a desired labeling $\\mathcal{L}^\\ast$ of $T^\\ast$.\n\nWe verify that $\\mathcal{L}^\\ast$ gives a correct labeling of $T^\\ast$. Clearly, the size-$\\Delta$ multiset that labels each vertex $v \\in V$ is in $\\mathcal{V}$, by the correctness of $\\mathcal{L}$ and $\\mathcal{L}''$. Consider any edge $e=\\{u,v\\}$ in $T^\\ast$. Similarly, if $\\{u,v\\} \\subseteq V \\setminus U$ or $\\{u,v\\} \\cap (V \\setminus U) = \\emptyset$, then the size-$2$ multiset that labels $e$ is in $\\mathcal{E}$, by the correctness of $\\mathcal{L}$ and $\\mathcal{L}''$. For the case that $e=\\{u,v\\}$ connects a vertex $u \\notin V$ and a vertex $v \\in V \\setminus U$, we must have $u \\in S'$ by the lemma statement. Therefore, the label of the half edge $(u,e)$ is the same in both $\\mathcal{L}$ and $\\mathcal{L}^\\ast$ by our choice of $\\mathcal{L}''$. Thus, the size-$2$ multiset that labels $e$ is in $\\mathcal{E}$, by the correctness of $\\mathcal{L}$.\n\nWe verify that the \nvirtual half-edge labels surrounding each pole in $S$ are the same in both $\\mathcal{L}$ and $\\mathcal{L}^\\ast$. Consider a pole $v \\in S$. If $v \\in V \\setminus U$, then the labeling of $v$ is clearly the same in both $\\mathcal{L}$ and $\\mathcal{L}^\\ast$.\nIf $v \\notin V\\setminus U$, then the condition $S \\cap U \\subseteq S'$ in the statement implies that $v \\in S'$. In this case, the way we pick $\\mathcal{L}''$ ensures that the \nvirtual half-edge labels surrounding $v \\in S'$ are the same in both $\\mathcal{L}$ and $\\mathcal{L}^\\ast$.\n\\end{proof}\n\nIn view of the proof of \\cref{lem:replace}, as long as the conditions in \\cref{lem:replace} are met, we are able to abstract out a subtree by its equivalence class when reasoning about correct labelings of a tree. This observation will be applied repeatedly in the subsequent discussion.\n\n\n\\paragraph{A Pumping Lemma} We will prove a pumping lemma of bipolar trees using \\cref{lem:replace}. \nSuppose $T_i$ is a tree with a root $r_i$ for each $1 \\leq i \\leq k$, then $H=(T_1, T_2, \\ldots, T_k)$ denotes the bipolar tree resulting from concatenating the roots $r_1, r_2, \\ldots, r_k$ into a path $(r_1, r_2, \\ldots, r_k)$ and setting the two poles of $H$ by $s = r_1$ and $t = r_k$.\nA simple consequence of \\cref{lem:replace} is that $\\texttt{type}(H)$ is determined by $\\textsf{Class}_1(T_1), \\textsf{Class}_1(T_2), \\ldots, \\textsf{Class}_1(T_k)$.\nWe have the following pumping lemma. \n\n\\begin{lemma}[Pumping lemma]\\label{lem:pump}\nThere exists a finite number $ {\\ell_{\\operatorname{pump}}} > 0$ such that as long as $k \\geq {\\ell_{\\operatorname{pump}}} $, any bipolar tree $H=(T_1, T_2, \\ldots, T_k)$ can be decomposed into $H = X \\circ Y \\circ Z$ with $0 < |Y| < k$ so that $X \\circ Y^i \\circ Z$ is in the same equivalence class as $H$ for each $i \\geq 0$. \n\\end{lemma}\n\\begin{proof}\nSet $ {\\ell_{\\operatorname{pump}}} $ to be the number of equivalence classes for bipolar trees plus one. By the pigeon hole principle, there exist $1 \\leq a < b \\leq k$ such that $(T_1, T_2, \\ldots, T_a)$ and $(T_1, T_2, \\ldots, T_b)$ are in the same equivalence class. Set $X = (T_1, T_2, \\ldots, T_a)$, $Y = (T_{a+1}, T_{a+2}, \\ldots, T_b)$, and $Z = (T_{b+1}, T_{b+2}, \\ldots, T_k)$. As we already know that $\\texttt{type}(X) = \\texttt{type}(X\\circ Y)$, \\cref{lem:replace} implies that \n$\\texttt{type}(X\\circ Y) = \\texttt{type}(X^2\\circ Y)$ by replacing $X$ by $X\\circ Y$ in the bipolar tree $X\\circ Y$.\nSimilarly, $\\texttt{type}(X \\circ Y^i)$ is the same for for each $i \\geq 0$. Applying \\cref{lem:replace} again to replace $X \\circ Y$ by $X \\circ Y^i$ in $H = X \\circ Y \\circ Z$, we conclude that $X \\circ Y^i \\circ Z$ is in the same equivalence class as $H$ for each $i \\geq 0$. \n\\end{proof}\n\n\n\n\\paragraph{A Procedure for Gluing Trees} Suppose that we have a set of rooted trees $W$. We devise a procedure that generates a new set of rooted trees by gluing the rooted trees in $W$ together. \nThis procedure is parameterized by a labeling function $f$. Consider a bipolar tree \\[H = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)\\] \nwhere $T^m$ is formed by attaching the roots $r_1, r_2, \\ldots, r_{\\Delta - 2}$ of the rooted trees $T_1^m, T_2^m, \\ldots, T_{\\Delta -2}^m$ to the root $r^m$ of $T^m$.\n\n\nThe labeling function $f$ assigns the half-edge labels surrounding $r^m$ based on \n\\[\\texttt{type}(H^l), \\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m), \\texttt{type}(H^r)\\]\nwhere $H^l = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l)$ and $H^r = (T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)$.\n\nWe write $T^{m}_\\ast$ to denote the result of applying $f$ to label the root $r^m$ of $T^{m}$ in the bipolar tree $H$, and we write \\[H_\\ast = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T_\\ast^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)\\] to denote the result of applying $f$ to $H$. We make the following observation.\n\n\\begin{lemma}[Property of $f$]\\label{lem:labeling_function}\nThe two equivalence classes\n$\\textsf{Class}_1(T_\\ast^m)$ and $\\texttt{type}(H_\\ast)$ are determined by \\[\\texttt{type}(H^l), \\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m), \\texttt{type}(H^r),\\] and the labeling function $f$.\n\\end{lemma}\n\\begin{proof}\nBy \\cref{lem:replace}, once the half-edge labelings of the root $r^m$ of $T^{m}$ is fixed, $\\textsf{Class}_1(T_\\ast^m)$ is determined by $\\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m)$. \nTherefore, indeed $\\textsf{Class}_1(T_\\ast^m)$ is determined by\n\\[\\texttt{type}(H^l), \\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m), \\texttt{type}(H^r),\\] and the labeling function $f$.\nSimilarly, applying \\cref{lem:replace} to the decomposition of $H_\\ast$ into $H^l$, $T_\\ast^m$, and $H^r$, we infer that $\\texttt{type}(H_\\ast)$ depends only on $\\texttt{type}(H^l)$, $\\textsf{Class}_1(T_\\ast^m)$, and $\\texttt{type}(H^r)$.\n\\end{proof}\n\nThe three sets of trees $X_f(W)$, $Y_f(W)$, and $Z_f(W)$ are constructed as follows.\n\n\\begin{itemize}\n \\item $X_f(W)$ is the set of all rooted trees resulting from appending $\\Delta-2$ arbitrary rooted trees $T_1, T_2, \\ldots, T_{\\Delta-2}$ in $W$ to a new root vertex $r$.\n \\item $Y_f(W)$ is the set of bipolar trees constructed as follows. For each choice of $2 {\\ell_{\\operatorname{pump}}} +1$ rooted trees $T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r$ from $X_f(W)$, concatenate them into a bipolar tree \n \\[H = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r),\\] let $H_\\ast$ be the result of applying the labeling function $f$ to $H$, and then add $H_\\ast$ to $Y_f(W)$.\n \\item $Z_f(W)$ is the set of rooted trees constructed as follows. For each \\[H_\\ast = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T_\\ast^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r) \\in Y_f(W),\\] add $(T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T_\\ast^m)$ to $Z_f(W)$, where we set the root of $T_1^l$ as the root, and add $(T_\\ast^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)$ to $Z_f(W)$, where we set the root of ${T_{ {\\ell_{\\operatorname{pump}}} }^r}$ as the root. \n\\end{itemize}\n\nWe write $\\textsf{Class}_1(S) = \\bigcup_{T \\in S} \\{\\textsf{Class}_1(T)\\}$ and $\\texttt{type}(S) = \\bigcup_{H \\in S} \\{\\texttt{type}(H)\\}$, and \nwe make the following observation.\n\n\\begin{lemma}[Property of $X_f(W)$, $Y_f(W)$, and $Z_f(W)$]\\label{lem:sets_of_trees}\nThe three sets of equivalence classes $\\textsf{Class}_1(X_f(W))$, $\\texttt{type}(Y_f(W))$, and $\\textsf{Class}_1(Z_f(W))$ depend only on $\\textsf{Class}_1(W)$ and the labeling function $f$.\n\\end{lemma}\n\\begin{proof}\nThis is a simple consequence of \\cref{lem:replace,lem:labeling_function}.\n\\end{proof}\n\n\\paragraph{A Fixed Point $W^\\ast$}\nGiven a fixed labeling function $f$,\nwe want to find a set of rooted trees $W^\\ast$ that is a \\emph{fixed point} for the procedure $Z_f$ in the sense that \\[\\textsf{Class}_1(Z_f(W^\\ast)) = \\textsf{Class}_1(W^\\ast).\\] \nTo find such a set $W^\\ast$, we construct a two-dimensional array of rooted trees $\\{W_{i,j}\\}$, as follows.\n\n\\begin{itemize}\n \\item For the base case, $W_{1,1}$ consists of only the one-vertex rooted tree.\n \\item Given that $W_{i,j}$ as been constructed, we define $W_{i,j+1} = Z_f(W_{i,j})$.\n \\item Given that $W_{i,j}$ for all positive integers $j$ have been constructed, $W_{i+1,1}$ is defined as follows. Pick $b_i$ as the smallest index such that $\\textsf{Class}_1(W_{i, b_i}) = \\textsf{Class}_1(W_{i, a_i})$ for some $1 \\leq a_i < b_i$. By the pigeon hole principle, the index $b_i$ exists, and it is upper bounded by $2^C$, where $C$ is the number of equivalence classes for rooted trees. We set \n \\[W_{i+1,1} = W_{i, a_i} \\cup W_{i, a_i + 1} \\cup \\cdots \\cup W_{i, b_i - 1}.\\]\n\\end{itemize}\n\nWe show that the sequence $\\textsf{Class}_1(W_{i,a_i}), \\textsf{Class}_1(W_{i,a_i+1}), \\textsf{Class}_1(W_{i,a_i+2}), \\ldots$ is periodic with a period $c_i = b_i - a_i$.\n\n\\begin{lemma}\\label{lem:W-aux-1}\nFor any $i \\geq 1$ and for any $j \\geq a_i$, we have $\\textsf{Class}_1(W_{i,j}) = \\textsf{Class}_1(W_{i,j+c_i})$, where $c_i = b_i - a_i$.\n\\end{lemma}\n\\begin{proof}\nBy \\cref{lem:sets_of_trees}, $\\textsf{Class}_1(W_{i,j})$ depends only on $\\textsf{Class}_1(W_{i,j-1})$. Hence the lemma follows from the fact that $\\textsf{Class}_1(W_{i, b_i}) = \\textsf{Class}_1(W_{i, a_i})$.\n\\end{proof}\n\nNext, we show that $\\textsf{Class}_1(W_{2,1}) \\subseteq \\textsf{Class}_1(W_{3, 1}) \\subseteq \\textsf{Class}_1(W_{4, 1}) \\subseteq \\cdots$.\n\n\\begin{lemma}\\label{lem:W-aux-2}\nFor any $i \\geq 2$, we have $\\textsf{Class}_1(W_{i,1}) \\subseteq \\textsf{Class}_1(W_{i,j})$ for each $j > 1$, and so $\\textsf{Class}_1(W_{i,1}) \\subseteq \\textsf{Class}_1(W_{i+1, 1})$.\n\\end{lemma}\n\\begin{proof}\nSince $W_{i,1} = \\bigcup_{a_{i-1} \\leq l \\leq b_{i-1} - 1} W_{i-1, l}$, we have \\[\\bigcup_{a_{i-1}+j-1 \\leq l \\leq b_{i-1} +j} W_{i-1, l} \\subseteq W_{i,j}\\] according to the procedure of constructing $W_{i,j}$. By \\cref{lem:W-aux-1}, \\[ \\textsf{Class}_1\\left(\\bigcup_{a_{i-1} \\leq l \\leq b_{i-1} - 1} W_{i-1, l}\\right) = \\textsf{Class}_1\\left(\\bigcup_{a_{i-1}+j-1 \\leq l \\leq b_{i-1} +j} W_{i-1, l}\\right)\\] for all $j \\geq 1$, and so we have $\\textsf{Class}_1(W_{i,1}) \\subseteq \\textsf{Class}_1(W_{i,j})$ for each $j > 1$. The claim $\\textsf{Class}_1(W_{i,1}) \\subseteq \\textsf{Class}_1(W_{i+1, 1})$ follows from the fact that $W_{i+1,1} = \\bigcup_{a_{i} \\leq l \\leq b_{i} - 1} W_{i, l}$.\n\\end{proof}\n\nSet $i^\\ast$ to be the smallest index $i \\geq 2$ such that $\\textsf{Class}_1(W_{i,1}) = \\textsf{Class}_1(W_{i+1, 1})$. By the pigeon hole principle and \\cref{lem:W-aux-2}, the index $i^\\ast$ exists, and it is upper bounded by $C$, the number of equivalence classes for rooted trees. We set \\[W^\\ast = W_{i^\\ast, 1}.\\] \n\nThe following lemma shows that $\\textsf{Class}_1(Z_f(W^\\ast)) = \\textsf{Class}_1(W^\\ast)$, as needed.\n\n\\begin{lemma}\\label{lem:W-aux-3}\nFor any $i \\geq 2$, if $\\textsf{Class}_1(W_{i,1}) = \\textsf{Class}_1(W_{i+1, 1})$, then $\\textsf{Class}_1(W_{i,j})$ is the same for all $j \\geq 1$. \n\\end{lemma}\n\\begin{proof}\nBy \\cref{lem:W-aux-1} and the way we construct $W_{i+1, 1}$, we have \\[\\textsf{Class}_1(W_{i, j}) \\subseteq \\textsf{Class}_1(W_{i+1, 1}) \\ \\ \\text{for each} \\ \\ j \\geq a_i.\\]\nBy \\cref{lem:W-aux-2}, we have \\[\\textsf{Class}_1(W_{i, 1}) \\subseteq \\textsf{Class}_1(W_{i, j}) \\ \\ \\text{for each} \\ \\ j > 1.\\]\nTherefore, $\\textsf{Class}_1(W_{i,1}) = \\textsf{Class}_1(W_{i+1, 1})$ implies that $\\textsf{Class}_1(W_{i, j}) = \\textsf{Class}_1(W_{i,1})$ for each $j \\geq a_i$. Hence we must have $\\textsf{Class}_1(Z_f(W_{i,1})) = \\textsf{Class}_1(W_{i,1})$, and so $\\textsf{Class}_1(W_{i,j})$ is the same for all $j \\geq 1$. \n\\end{proof}\n\nWe remark that simply selecting $W^\\ast$ to be any set such that $\\textsf{Class}_1(Z_f(W^\\ast)) = \\textsf{Class}_1(W^\\ast)$ is not enough for our purpose. As we will later see, it is crucial that the set $W^\\ast$ is constructed by iteratively applying the function $Z_f$ and taking the union of previously constructed sets.\n\n\n\n\n\n\n\n\\paragraph{A Sequence of Sets of Trees}\nWe define $W_1^\\ast = W^\\ast$ and $W_i^\\ast = Z_f(W_{i-1}^\\ast)$ for each $1 < i \\leq L$, where $L$ is some sufficiently large number to be determined.\n\nThe way we choose $W^\\ast$ guarantees that $\\textsf{Class}_1(W_i^\\ast) = \\textsf{Class}_1(W^\\ast)$ for all $1 \\leq i \\leq L$. \nFor convenience, we write $X_i^\\ast = X_f(W_i^\\ast)$ and $Y_i^\\ast = Y_f(W_i^\\ast)$.\nSimilarly, $\\textsf{Class}_1(X_i^\\ast)$ is the same for all $1 \\leq i \\leq L$ and $\\texttt{type}(Y_i^\\ast)$ is the same for all $1 \\leq i \\leq L$.\n\n\nOur analysis will rely on the assumption that all rooted trees in $W^\\ast$ admit correct labelings. Whether this is true depends only on $\\textsf{Class}_1(W^\\ast)$, which depends only on the labeling function $f$. We say that $f$ is \\emph{feasible} if it leads to a set $W^\\ast$ where all the rooted trees therein admit correct labelings. The proof that a feasible labeling function $f$ exists is deferred.\n\nThe assumption that $f$ is feasible implies that all trees in $W_i^\\ast$, $X_i^\\ast$, and $Y_i^\\ast$, for all $1 \\leq i \\leq L$, admit correct labelings. All rooted trees in $X_i^\\ast$ admit correct labelings because they are subtrees of the rooted trees in $W_{i+1}^\\ast$. A correct labeling of any bipolar tree $H \\in Y_i^\\ast$ can be obtained by combining any correct labelings of the two rooted trees in $W_{i+1}^\\ast$ resulting from $H$.\n\n\n\n\n\\paragraph{Layers of Vertices}\nWe assign a layer number $\\lambda(v)$ to each vertex $v$ in a tree based on the step that $v$ is introduced in the construction \\[W_1^\\ast \\rightarrow W_2^\\ast \\rightarrow \\cdots \\rightarrow W_L^\\ast.\\]\n\nIf a vertex $v$ is introduced as the root vertex of a tree in $X_i^\\ast = X_f(W_i^\\ast)$, then we say that the layer number of $v$ is $\\lambda(v) = i \\in \\{1,2, \\ldots, L\\}$. A vertex $v$ has $\\lambda(v) = 0$ if it belongs to a tree in $W_1^\\ast$.\n\n\nFor any vertex $v$ with $\\lambda(v) = i$ in a tree $T \\in X_j^\\ast$ with $i \\leq j$, we write $T_v$ to denote the subtree of $T$ such that $T_v \\in X_i^\\ast$ where $v$ is the root of $T_v$. \n\n\n \nWe construct a sequence of sets $R_1, R_2, \\ldots, R_L$ as follows.\nWe go over all rooted trees $T \\in X_L^\\ast$, all possible correct labeling $\\mathcal{L}$ of $T$, and all vertices $v$ in $T$ with $\\lambda(v) = i \\in \\{1,2, \\ldots, L\\}$. Suppose that the $\\Delta-2$ rooted trees in the construction of $T_v \\in X_i^\\ast = X_f(W_i^\\ast)$ are $T_1$, $T_2$, $\\ldots$, $T_{\\Delta-2}$, and let $r_i$ be the root of $T_i$. Consider the following parameters.\n\\begin{itemize}\n \\item $c_i = \\textsf{Class}_1(T_i)$.\n \\item $\\mathtt{a}_i$ is the real half-edge label of $v$ in $T_v$ for the edge $\\{v, r_i\\}$, under the correct labeling $\\mathcal{L}$ of $T$ restricted to $T_v$.\n \\item $\\mathcal{I}$ is the size-$2$ multiset of virtual half-edge labels of $v$ in $T_v$, under the correct labeling $\\mathcal{L}$ of $T$ restricted to $T_v$.\n\\end{itemize}\nThen we add \n$(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I})$ to $R_i$.\n\n\n\n\n\\begin{lemma}\\label{lem:Rsets1}\nFor each $1 \\leq i < L$, $R_i$ is determined by $R_{i+1}$.\n\\end{lemma}\n\\begin{proof}\nWe consider the following alternative way of constructing $R_i$ from $R_{i+1}$.\nEach rooted tree $T'$ in $W_{i+1}^\\ast = Z_f(W_i^\\ast)$ can be described as follows.\n\\begin{itemize}\n \\item Start with a path $(r_1, r_2, \\ldots, r_{ {\\ell_{\\operatorname{pump}}} + 1})$, where $r_1$ is the root of $T'$.\n \\item For each $1 \\leq j \\leq {\\ell_{\\operatorname{pump}}} +1$, append $\\Delta-2$ rooted trees $T_{j,1}, T_{j,2}, \\ldots, T_{j, \\Delta-2} \\in W_i^\\ast$ to $r_j$.\n \\item Assign the labels to the half edges surrounding $r_{ {\\ell_{\\operatorname{pump}}} + 1}$ according to the labeling function $f$.\n\\end{itemize}\n\n Now, consider the function $\\phi$ that maps each equivalence class $c$ for rooted trees to a subset of $\\Sigma$ defined as follows:\n$\\mathtt{a} \\in \\phi(c)$ if there exist \\[(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I}) \\in R_{i+1}\\] and $1 \\leq j \\leq \\Delta -2$ such that $c = c_j$ and $\\mathtt{a} = \\mathtt{a}_j$.\n\nWe go over all possible $T' \\in W_{i+1}^\\ast = Z_f(W_i^\\ast)$.\nNote that the root $r$ of $T'$ has exactly one virtual half edge. For each $\\mathtt{b} \\in \\Sigma$ such that $\\{\\mathtt{a}, \\mathtt{b}\\} \\in \\mathcal{E}$ for some $\\mathtt{a} \\in \\phi(\\textsf{Class}_1(T'))$, we go over all possible correct labelings $\\mathcal{L}$ of $T'$ where the virtual half edge of $r$ is labeled $\\mathtt{b}$. For each $1 \\leq j \\leq {\\ell_{\\operatorname{pump}}} +1$, consider the following parameters.\n\\begin{itemize}\n \\item $c_l = \\textsf{Class}_1(T_{j,l})$.\n \\item $\\mathtt{a}_l$ is the half-edge label of $r_j$ for the edge $\\{r_j, r_{j,l}\\}$.\n \\item $\\mathcal{I}$ is the size-$2$ multiset of the remaining two half-edge labels of $v$.\n\\end{itemize}\nThen we add \n$(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I})$ to $R_i$. \n\nThis construction of $R_i$ is equivalent to the original construction of $R_i$ because a correct labeling of $T' \\in W_i^\\ast$ can be extended to a correct labeling of a tree $T \\in X_L^\\ast$ that contains $T'$ as a subtree if and only if the virtual half-edge label of the root $r$ of $T'$ is $\\mathtt{b} \\in \\Sigma$ such that $\\{\\mathtt{a}, \\mathtt{b}\\} \\in \\mathcal{E}$ for some $\\mathtt{a} \\in \\phi(\\textsf{Class}_1(T'))$.\n\n\nIt is clear that this construction of $R_i$ only depends on $\\textsf{Class}_1(W_i^\\ast)$, the labeling function $f$, and the function $\\phi$, which depends only on $R_{i+1}$. Since the labeling function $f$ is fixed and $\\textsf{Class}_1(W_i^\\ast)$ is the same for all $i$, we conclude that $R_{i}$ depends only on $R_{i+1}$.\n\\end{proof}\n\n\\begin{lemma}\\label{lem:Rsets2}\nWe have $R_1 \\subseteq R_2 \\subseteq \\cdots \\subseteq R_L$.\n\\end{lemma}\n\\begin{proof}\nFor the base case, we show that $R_{L-1} \\subseteq R_L$. In fact, our proof will show that $R_{i} \\subseteq R_L$ for each $1 \\leq i < L$. Consider any \n\\[(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I}) \\in R_i.\\] \nThen there is a rooted tree $T \\in X_i^\\ast$ that is formed by attaching $\\Delta-2$ rooted trees of equivalence classes $c_1, c_2, \\ldots, c_{\\Delta-2}$ to the root vertex so that if we label the half edges surrounding the root vertex according to $\\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}$, and $\\mathcal{I}$, then this partial labeling can be completed into a correct labeling of $T$.\n\nBecause $\\textsf{Class}_1(W_i^\\ast) = \\textsf{Class}_1(W_L^\\ast)$, there is also a rooted tree $T' \\in X_L^\\ast$ that is formed by attaching $\\Delta-2$ rooted trees of equivalence classes $c_1, c_2, \\ldots, c_{\\Delta-2}$ to the root vertex. Therefore, if we label the root vertex of $T'$ in the same way as we do for $T$, then this partial labeling can also be completed into a correct labeling of $T'$. Hence we must have \n\\[(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I}) \\in R_L.\\] \n\n\nNow, suppose that we already have $R_i \\subseteq R_{i+1}$ for some $1 < i < L$. We will show that $R_{i-1} \\subseteq R_{i}$. \nDenote by $\\phi_{i}$ and $\\phi_{i+1}$ the function $\\phi$ in \\cref{lem:Rsets1} constructed from $R_{i}$ and $R_{i+1}$. \nWe have $\\phi_i(c) \\subseteq \\phi_{i+1}(c)$ for each equivalence class $c$, because $R_i \\subseteq R_{i+1}$. Therefore, in view of the alternative construction described in the proof of \\cref{lem:Rsets1}, we have $R_{i-1} \\subseteq R_{i}$.\n\\end{proof}\n\n\n\n\\paragraph{The Set of Vertex Configurations $\\mathcal{V}'$}\nBy \\cref{lem:Rsets1,lem:Rsets2}, if we pick $L$ to be sufficient large, we can have $R_1 = R_2 = R_3$.\nMore specifically, if we pick \n\\[L \\geq C^{\\Delta-2} \\cdot \\binom{|\\Sigma|+1}{|\\Sigma|-1} + 3,\\]\nthen there exists an index $3 \\leq i \\leq L$ such that $R_i = R_{i-1}$, implying that $R_1 = R_2 = R_3$. Here $C$ is the number of equivalence classes for rooted trees and $\\binom{|\\Sigma|+1}{|\\Sigma|-1}$ is the number of size-$2$ multisets of elements from $\\Sigma$.\n\n\n\nThe set $\\mathcal{V}'$ is defined by including all size-$\\Delta$ multisets $\\{\\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}\\} \\cup \\mathcal{I}$ such that \\[(c_1, c_2, \\ldots, c_{\\Delta-2}, \\mathtt{a}_1, \\mathtt{a}_2, \\ldots, \\mathtt{a}_{\\Delta-2}, \\mathcal{I}) \\in R_{1}\\] for some $c_1, c_2, \\ldots, c_{\\Delta-2}$.\n\n\n\n\\paragraph{The Set $\\mathcal{V}'$ is $\\ell$-full}\nTo show that the set $\\mathcal{V}'$ is $\\ell$-full, we consider the subset $\\mathcal{V}^\\ast \\subseteq \\mathcal{V}'$ defined by the set of vertex configurations used by the labeling function $f$ in the construction of $Y_f(W_i^\\ast)$. The definition of $\\mathcal{V}^\\ast$ is invariant of $i$ as $\\textsf{Class}_1(W_i^\\ast)$ is the same for all $i$. Clearly, for each $x \\in \\mathcal{V}^\\ast$, there is a rooted tree $T \\in W_2^\\ast$ where the root of a subtree $T' \\in X_1^\\ast$ of $T$ has its size-$\\Delta$ multiset of half-edge labels fixed to be $x$ by $f$, and so $\\mathcal{V}^\\ast \\subseteq \\mathcal{V}'$.\n\nFor notational simplicity, we write $x \\overset{k}{\\leftrightarrow} x'$ if there is a correct labeling of a $k$-vertex path $(v_1, v_2, \\ldots, v_k)$ using only vertex configurations in $\\mathcal{V}'$ so that the vertex configuration of $v_1$ is $x$ and the vertex configuration of $v_k$ is $x'$. If it is further required that the half-edge label of $v_1$ for the edge $\\{v_1, v_2\\}$ is $\\mathtt{a} \\in x$, then we write $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} x'$. The notation $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{a}')$ is defined similarly.\n\n\n\\begin{lemma}\\label{lem:find_path_1}\nFor any $x \\in \\mathcal{V}' \\setminus \\mathcal{V}^\\ast$ and $\\mathtt{a} \\in x$, there exist a vertex configuration $x' \\in \\mathcal{V}^\\ast$ and a number $2 \\leq k \\leq 2 {\\ell_{\\operatorname{pump}}} + 1$ such that $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} x'$.\n\\end{lemma}\n\\begin{proof}\nLet $T \\in W_L^\\ast$ be chosen so that there is a correct labeling $\\mathcal{L}$ where $x$ is a vertex configuration of some vertex $v$ with $\\lambda(v) = 2$.\n\nTo prove the lemma, it suffices to show that for each of the $\\Delta$ neighbors $u$ of $v$, it is possible to find a path $P = (v,u, \\ldots, w)$ meeting the following conditions.\n\\begin{itemize}\n \\item $w$ is a vertex whose half-edge labels have been fixed by $f$. This ensures that the vertex configuration of $w$ is in $\\mathcal{V}^\\ast$.\n \\item All vertices in $P$ are within layers 1,2, and 3. This ensures that the vertex configuration of all vertices in $P$ are in $\\mathcal{V}'$.\n \\item The number of vertices $k$ in $P$ satisfies $2 \\leq k \\leq 2 {\\ell_{\\operatorname{pump}}} + 1$. \n\\end{itemize}\n\nWe divide the proof into three cases. Refer to \\cref{fig:ell_full_construction1} for an illustration, where squares are vertices whose half-edge labels have been fixed by $f$.\n\n\\paragraph{Case 1} Consider the subtree $T_v \\in X_2^\\ast$ of $T$ whose root is $v$. In view of the construction of the set $X_f(W_2^\\ast)$, $v$ has $\\Delta-2$ children $u_1, u_2, \\ldots, u_{\\Delta-2}$ in $T_v$, where the subtree $T_i$ rooted at $u_i$ is a rooted tree in $W_2^\\ast$. \n\nFor each $1 \\leq i \\leq \\Delta-2$, according to the structure of the trees in the set $W_2^\\ast = Z_f(W_1^\\ast)$, there is a path $(u_i = w_1, w_2, \\ldots, w_{ {\\ell_{\\operatorname{pump}}} +1})$ in $T_i$ containing only layer-1 vertices, where the half-edge labels of $w_{ {\\ell_{\\operatorname{pump}}} +1}$ have been fixed by $f$. Hence $P = (v, u_i = w_1, w_2, \\ldots, w_{ {\\ell_{\\operatorname{pump}}} +1})$ is a desired path with $k = {\\ell_{\\operatorname{pump}}} + 2$ vertices.\n \n \n \n \n \n\\paragraph{Case 2} Consider the subtree $T' \\in W_3^\\ast = Z_f(W_2^\\ast)$ that contains $v$ in $T$.\nSimilarly, according to the structure of the trees in the set $W_3^\\ast = Z_f(W_2^\\ast)$, there is a path $(r = w_1', w_2', \\ldots, w_{ {\\ell_{\\operatorname{pump}}} +1}')$ in $T'$ containing only layer-2 vertices so that $r$ is the root of $T'$, $v = w_i'$ for some $1 \\leq i' \\leq {\\ell_{\\operatorname{pump}}} +1$, and the half-edge labels of $w_{ {\\ell_{\\operatorname{pump}}} +1}'$ have been fixed by $f$.\nSince $x \\in \\mathcal{V}' \\setminus \\mathcal{V}^\\ast$, we have $v \\neq w_{ {\\ell_{\\operatorname{pump}}} +1}'$. Hence $P = (v = w_i', w_{i+1}', \\ldots, w_{ {\\ell_{\\operatorname{pump}}} +1}')$ is a desired path with $2 \\leq k \\leq {\\ell_{\\operatorname{pump}}} + 1$ vertices.\n\n \n \n \n\n\\paragraph{Case 3} There is only one remaining neighbor of $v$ to consider. In view of the construction of $X_f(W_3^\\ast)$, there is a layer-3 vertex $v'$ adjacent to the vertex $r$, the root of the tree $T' \\in W_3^\\ast$ considered in the previous case. If the half-edge labels of $v'$ have been fixed by $f$, then $P = (v = w_i', w_{i-1}', \\ldots, w_{1}'=r, v')$ is a desired path.\nOtherwise, similar to the analysis in the previous case, we can find a path $P' = (v', \\ldots, w)$ connecting $v'$ to a vertex $w$ whose half-edge labels have been fixed by $f$. All vertices in $P'$ are of layer-3, and the number of vertices in $P'$ is within $[2, {\\ell_{\\operatorname{pump}}} + 1]$. \nCombining $P'$ with the path $(v = w_i', w_{i-1}', \\ldots, w_{1}'=r, v')$, we obtain the desired path $P$ whose number of vertices satisfies $2 \\leq k \\leq 2 {\\ell_{\\operatorname{pump}}} + 1$.\n \n \n\\end{proof}\n\n\n\\begin{figure}[h!]\n \\centering\n \\includegraphics[width= .75\\textwidth]{fig\/y2.pdf}\n \\caption{An illustration for the proof of \\cref{lem:find_path_1}.}\n \\label{fig:ell_full_construction1}\n\\end{figure}\n\nFor \\cref{lem:find_path_2}, note that if $\\mathtt{a}$ appears more than once in the multiset $x$, then we still have $\\mathtt{a} \\in x \\setminus \\{\\mathtt{a}\\}$.\n\n\n\\begin{lemma}\\label{lem:find_path_2}\nFor any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$, $\\mathtt{a}' \\in x' \\in \\mathcal{V}^\\ast$, and $0 \\leq t \\leq {\\ell_{\\operatorname{pump}}} - 1$, we have $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ for some $\\mathtt{b} \\in x \\setminus \\{\\mathtt{a}\\}$ and $\\mathtt{b}' \\in x' \\setminus \\{\\mathtt{a}'\\}$\n with $k = 2 {\\ell_{\\operatorname{pump}}} + 4 + t$.\n\\end{lemma}\n\\begin{proof}\nFor any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$, we can find a rooted tree $T_{x, \\mathtt{a}} \\in W_2^\\ast$ such that the path $(v_1$, $v_2$, $\\ldots$, $v_{ {\\ell_{\\operatorname{pump}}} +1})$ of the layer-1 vertices in $T_{x, \\mathtt{a}}$ where $v_1 = r$ is the root of $T_{x, \\mathtt{a}}$ satisfies the property that the half-edge labels of $v_{ {\\ell_{\\operatorname{pump}}} +1}$ have been fixed to $x$ by $f$ where the half-edge label for $\\{v_{ {\\ell_{\\operatorname{pump}}} }, v_{ {\\ell_{\\operatorname{pump}}} +1}\\}$ is in $x \\setminus \\{\\mathtt{a}\\}$.\n\nNow, observe that for any choice of $(\\Delta-2)( {\\ell_{\\operatorname{pump}}} +1)$ rooted trees \\[\\{T_{i,j}\\}_{1 \\leq i \\leq {\\ell_{\\operatorname{pump}}} +1, 1 \\leq j \\leq \\Delta-2}\\] in $W_2^\\ast$, there is a rooted tree $T' \\in W_3^\\ast = Z_f(W_2^\\ast)$ formed by appending $T_{i, 1}$, $T_{i,2}$, $\\ldots$, $T_{i, \\Delta-2}$ to $u_i$ in the path $(u_1$, $u_2$, $\\ldots$, $u_{ {\\ell_{\\operatorname{pump}}} +1})$ and fixing the half-edge labeling of $u_{ {\\ell_{\\operatorname{pump}}} +1}$ by $f$. All vertices in $(u_1, u_2, \\ldots, u_{ {\\ell_{\\operatorname{pump}}} +1})$ are of layer-2.\n\nWe choose any $T' \\in W_2^\\ast$ with $T_{i,j} = T_{x, \\mathtt{a}}$ and $T_{i',j'} = T_{x', \\mathtt{a}'}$ such that $i' - i = t + 1$.\nThe possible range of $t$ is $[0, {\\ell_{\\operatorname{pump}}} - 1]$.\nConsider any $T \\in W_L^\\ast$ that contains $T'$ as its subtree, and consider any correct labeling of $T$.\nThen the path $P$ resulting from concatenating the path $(v_{ {\\ell_{\\operatorname{pump}}} +1}, v_{ {\\ell_{\\operatorname{pump}}} }, \\ldots, v_1)$ in $T_{x, \\mathtt{a}}$, the path $(u_i, u_{i+1}, \\ldots, u_{i'})$, and the path $(v_1', v_2', \\ldots, v_{ {\\ell_{\\operatorname{pump}}} +1}')$ in $T_{x', \\mathtt{a}'}$\nshows that $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ \nfor $k = ( {\\ell_{\\operatorname{pump}}} +1) + (t+2) + ( {\\ell_{\\operatorname{pump}}} +1) = 2 {\\ell_{\\operatorname{pump}}} + 4 + t$. See \\cref{fig:ell_full_construction2} for an illustration.\n\nFor the rest of the proof, we show that the desired rooted tree $T_{x, \\mathtt{a}} \\in W_2^\\ast$ exists for any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$.\nFor any $x \\in \\mathcal{V}^\\ast$, we can find a bipolar tree\n\\[H_\\ast = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T_\\ast^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r) \\in Y_1^\\ast = Y_f(W_1^\\ast)\\]\nsuch that the vertex configuration of the root $r^m$ of $T_\\ast^m$ is fixed to be $x$ by the labeling function $f$.\nThen, for any $\\mathtt{a} \\in x$, at least one of the two rooted trees in $W_2^\\ast = Z_f(W_1^\\ast)$ resulting from cutting $H_\\ast$ satisfies the desired requirement. \n\\end{proof}\n\n\n\\begin{figure}[h!]\n \\centering\n \\includegraphics[width= .65\\textwidth]{fig\/y3.pdf}\n \\caption{An illustration for the proof of \\cref{lem:find_path_2}.}\n \\label{fig:ell_full_construction2}\n\\end{figure} \n\nIn the subsequent discussion, the length of a path refers to the number of edges in a path. In particular, the length of a $k$-vertex path is $k-1$.\n\nThe following lemma is proved by iteratively applying \\cref{lem:find_path_2} via intermediate vertex configurations $\\tilde{x} \\in \\mathcal{V}^\\ast$. \n\n\\begin{lemma}\\label{lem:find_path_3}\nFor any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$ and $\\mathtt{a}' \\in x' \\in \\mathcal{V}^\\ast$, we have $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ for some $\\mathtt{b} \\in x \\setminus \\{\\mathtt{a}\\}$ and $\\mathtt{b}' \\in x' \\setminus \\{\\mathtt{a}'\\}$\n for all $k \\geq 6 {\\ell_{\\operatorname{pump}}} + 10$.\n\\end{lemma}\n\\begin{proof}\nBy applying \\cref{lem:find_path_2} for three times,\nwe infer that \\cref{lem:find_path_2} also works for any path length $k - 1 = (k_1 - 1) + (k_2 - 1) + (k_3 - 1)$, where $k_i - 1 = 2 {\\ell_{\\operatorname{pump}}} + 3 + t_i$ for $0 \\leq t_i \\leq {\\ell_{\\operatorname{pump}}} -1$. \n\nSpecifically, we first apply \\cref{lem:find_path_2} to find a path realizing $(x,\\mathtt{b}) \\overset{k_1}{\\leftrightarrow} \\tilde{x}$ for some $\\tilde{x} \\in \\mathcal{V}^\\ast$, and for some $\\mathtt{b} \\in x \\setminus \\{\\mathtt{a}\\}$.\nLet $\\tilde{\\mathtt{a}} \\in \\tilde{x}$ be the real half-edge label used in the last vertex of the path. \nWe extend the length of this path by $k_2 - 1$ by attaching to it the path realizing $(\\tilde{x},\\tilde{\\mathtt{b}}) \\overset{k_2}{\\leftrightarrow} \\tilde{x}$, for some $\\tilde{\\mathtt{b}} \\in \\tilde{x} \\setminus \\{\\tilde{\\mathtt{a}}\\}$.\nFinally, let $\\tilde{\\mathtt{c}} \\in \\tilde{x}$ be the real half-edge label used in the last vertex of the current path. \nWe extend the length of the current path by $k_3$ by attaching to it the path realizing $(\\tilde{x},\\tilde{\\mathtt{d}}) \\overset{k_3}{\\leftrightarrow} (x', \\mathtt{b}')$, for some $\\tilde{\\mathtt{d}} \\in \\tilde{x} \\setminus \\{\\tilde{\\mathtt{x}}\\}$, and for some $\\mathtt{b}' \\in x' \\setminus \\{\\mathtt{a}'\\}$.\nThe resulting path realizes $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ for some $\\mathtt{b} \\in x \\setminus \\{\\mathtt{a}\\}$ and $\\mathtt{b}' \\in x' \\setminus \\{\\mathtt{a}'\\}$.\n\n\nTherefore, \\cref{lem:find_path_2} also works with path length of the form $k-1 = 6 {\\ell_{\\operatorname{pump}}} + 9 + t$, for any $t \\in [0, 3 {\\ell_{\\operatorname{pump}}} - 3]$. \n\nAny $k-1 \\geq 6 {\\ell_{\\operatorname{pump}}} + 9$ can be written as $k-1 = b(2 {\\ell_{\\operatorname{pump}}} +3) + (6 {\\ell_{\\operatorname{pump}}} + 9 + t)$, for some $t \\in [0, 3 {\\ell_{\\operatorname{pump}}} - 3]$ and some integer $b \\geq 0$.\nTherefore, similar to the above, we can find a path showing $(x,\\mathtt{b}) \\overset{k}{\\leftrightarrow} (x', \\mathtt{b}')$ by first applying \\cref{lem:find_path_2} with path length\n$2 {\\ell_{\\operatorname{pump}}} +3$ for $b$ times, and then applying the above variant of \\cref{lem:find_path_2} to extend the path length by $6 {\\ell_{\\operatorname{pump}}} + 9 + t$.\n\\end{proof}\n\nWe show that \\cref{lem:find_path_1,lem:find_path_3} imply that $\\mathcal{V}'$ is $\\ell$-full for some $\\ell$. \n\n\\begin{lemma}[$\\mathcal{V}'$ is $\\ell$-full] \\label{lem:find_path_4}\nThe set $\\mathcal{V}'$ is $\\ell$-full for $\\ell = 10 {\\ell_{\\operatorname{pump}}} + 10$. \n\\end{lemma}\n\\begin{proof}\nWe show that for any target path length $k-1 \\geq 10 {\\ell_{\\operatorname{pump}}} + 9$, and for any $\\mathtt{a} \\in x \\in \\mathcal{V}^\\ast$ and $\\mathtt{a}' \\in x' \\in \\mathcal{V}^\\ast$, we have $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} (x',\\mathtt{a}')$.\n\n\nBy \\cref{lem:find_path_1}, there exists a vertex configuration $\\tilde{x} \\in \\mathcal{V}^\\ast$ so that we can find a path $P$ realizing $(x,\\mathtt{a}) \\overset{\\ell_1}{\\leftrightarrow} \\tilde{x}$ for some path length $1 \\leq (\\ell_1-1) \\leq 2 {\\ell_{\\operatorname{pump}}} $.\nSimilarly, there exists a vertex configuration $\\tilde{x}' \\in \\mathcal{V}^\\ast$ so that we can find a path $P'$ realizing $(x',\\mathtt{a}') \\overset{\\ell_2}{\\leftrightarrow} \\tilde{x}'$ for some path length $1 \\leq (\\ell_2-1) \\leq 2 {\\ell_{\\operatorname{pump}}} $. \n\n\nLet $\\tilde{\\mathtt{a}}$ be the real half-edge label for $\\tilde{x}$ in $P$, and let $\\tilde{\\mathtt{a}}'$ be the real half-edge label for $\\tilde{x}'$ in $P'$.\nWe apply \\cref{lem:find_path_3} to find a path $\\tilde{P}$ realizing $(x,\\mathtt{b}) \\overset{\\tilde{\\ell}}{\\leftrightarrow} (x', \\mathtt{b}')$ for some $\\mathtt{b} \\in \\tilde{x} \\setminus \\{\\tilde{\\mathtt{a}}\\}$ and $\\mathtt{b}' \\in \\tilde{x}' \\setminus \\{\\tilde{\\mathtt{a}}'\\}$ with path length \n\\[\\tilde{\\ell}-1 = (k-1) - (\\ell_1-1) - (\\ell_2-1) \\geq 6 {\\ell_{\\operatorname{pump}}} + 9.\\] \nThe path formed by concatenating $P_1$, $\\tilde{P}$, and $P_2$ shows that $(x,\\mathtt{a}) \\overset{k}{\\leftrightarrow} (x',\\mathtt{a}')$.\n\\end{proof}\n\n\n\n\n\n\n \n\n\n\n\\paragraph{A Feasible Labeling Function $f$ exists} We show that a feasible labeling function $f$ exists given that $\\Pi$ can be solved in $\\mathsf{LOCAL}(n^{o(1)})$ rounds. We will construct a labeling function $f$ in such a way each equivalence class in $\\textsf{Class}_1(W^\\ast)$ contains only rooted trees that admit legal labeling. That is, for each $c \\in \\textsf{Class}_1(W^\\ast)$, the mapping $h$ associated with $c$ satisfies $h(X) = \\textsf{YES}$ for some $X$.\n\n\n\nWe will consider a series of modifications in the construction \\[W \\rightarrow X_f(W) \\rightarrow Y_f(W) \\rightarrow Z_f(W)\\] that do not alter the sets of equivalence classes in these sets, conditioning on the assumption that all trees that we have processed admit correct labelings. That is, whether $f$ is feasible is invariant of the modifications.\n\n\n\\paragraph{Applying the Pumping Lemma} Let $w > 0$ be some target length for the pumping lemma. In the construction of $Y_f(W)$, when we process a bipolar tree \n\\[H = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l, T^m, T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r),\\] we apply \\cref{lem:pump} to the two subtrees $H^l = (T_1^l, T_2^l, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^l)$ and $H^r = (T_1^r, T_2^r, \\ldots, T_{ {\\ell_{\\operatorname{pump}}} }^r)$ to obtain two new bipolar trees $H_+^l$ and $H_+^r$.\nThe $s$-$t$ path in the new trees $H_+^l$ and $H_+^r$ contains $w + x$ vertices, for some $0 \\leq x < {\\ell_{\\operatorname{pump}}} $. The equivalence classes do not change, that is, $\\texttt{type}(H^l) = \\texttt{type}(H_+^l)$ and $\\texttt{type}(H^r) = \\texttt{type}(H_+^r)$.\n\n\nWe replace $H^l$ by $H_+^l$ and replace $H^r$ by $H_+^r$ in the bipolar tree $H$.\nRecall that the outcome of applying the labeling function $f$ to the root $r$ of the rooted tree $T_{ {\\ell_{\\operatorname{pump}}} +1}$ depends only on \\[\\texttt{type}(H^l), \\textsf{Class}_1(T_{1}^m), \\textsf{Class}_1(T_{2}^m), \\ldots, \\textsf{Class}_1(T_{\\Delta - 2}^m), \\texttt{type}(H^r),\\] so applying the pumping lemma to $H^l$ and $H^r$ during the construction of $Y_f(W)$ does not alter $\\textsf{Class}_1(T_\\ast^m)$ and $\\texttt{type}(H^\\ast)$ for the resulting bipolar tree $H^\\ast$ and its middle rooted tree $T_\\ast^m$, by \\cref{lem:labeling_function}.\n\n\n\n\\paragraph{Reusing Previous Trees} During the construction of the fixed point $W^\\ast$, we remember all bipolar trees to which we have applied the feasible function $f$.\n\nDuring the construction of $Y_f(W)$. Suppose that we are about to process a bipolar tree $H$, and there is already some other bipolar tree $\\tilde{H}$ to which we have applied $f$ before so that \n$\\texttt{type}(H^l)=\\texttt{type}(\\tilde{H}^l)$, $\\textsf{Class}_1(T_{i}^m) = \\textsf{Class}_1(\\tilde{T}_{i}^m)$ for each $1 \\leq i \\leq \\Delta-2$, and $\\texttt{type}(H^r)=\\texttt{type}(\\tilde{H}^r)$.\nThen we replace $H$ by $\\tilde{H}$, and then we process $\\tilde{H}$ instead.\n\nBy \\cref{lem:labeling_function}, this modification does not alter $\\textsf{Class}_1(T_\\ast^m)$ and $\\texttt{type}(H^\\ast)$ for the resulting bipolar tree $H^\\ast$.\n\n\\paragraph{Not Cutting Bipolar Trees} We consider the following different construction of $Z_f(W)$ from $Y_f(W)$. For each $H^\\ast \\in Y_f(W)$, we simply add two copies of $H^\\ast$ to $Z_f(W)$, one of them has $r = s$ and the other one has $r = t$.\nThat is, we do not cut the bipolar tree $H^\\ast$, as in the original construction of $Z_f(W)$.\n\nIn general, this modification might alter $\\textsf{Class}_1(Z_f(W))$.\nHowever, it is not hard to see that if all trees in $Y_f(W)$ already admit correct labelings, then $\\textsf{Class}_1(Z_f(W))$ does not alter after the modification.\n\nSpecifically, let $T$ be a rooted tree in $Z_f(W)$ with the above modification. \nThen $T$ is identical to a bipolar tree $H^\\ast=(T_1, T_2, \\ldots, T_k) \\in Y_f(W)$, where there exists some $1 < i < k$ such that the root $r_i$ of $T_i$ has its half-edge labels fixed by $f$.\nSuppose that the root of $T$ is the root $r_1$ of $T_1$.\nIf we do not have the above modification, then the rooted tree $T'$ added to $Z_f(W)$ corresponding to $T$ is $(T_1, T_2, \\ldots, T_i)$, where the root of $T'$ is also $r_1$. \n\nGiven the assumption that all bipolar trees in $Y_f(W)$ admit correct labelings, it is not hard to see that $\\textsf{Class}_1(T') = \\textsf{Class}_1(T)$. For any correct labeling $\\mathcal{L}'$ of $T'$, we can extend the labeling to a correct labeling $\\mathcal{L}$ of $T$ by labeling the remaining vertices according to any arbitrary correct labeling of $H^\\ast$. This is possible because the labeling of $r_i$ has been fixed. For any correct labeling $\\mathcal{L}$ of $T$, we can obtain a correct labeling $\\mathcal{L}'$ of $T'$ by restricting $\\mathcal{L}$ to $T'$.\n\n\n\n\n\n\n\\paragraph{Simulating a $\\mathsf{LOCAL}$ Algorithm} We will show that there is a labeling function $f$ that makes all the rooted trees in $W^\\ast$ to admit correct solutions, where we apply the above three modifications in the construction of $W^\\ast$. In view of the above discussion, such a function $f$ is feasible.\n\nThe construction of $W^\\ast$ involves only a finite number $k$ of iterations of $Z_f$ applied to some previously constructed rooted trees. If we view the target length $w$ of the pumping lemma as a variable, then the size of a tree in $W^\\ast$ can be upper bounded by $O(w^{k})$. \n\n\nSuppose that we are given an arbitrary $n^{o(1)}$-round $\\mathsf{LOCAL}$ algorithm $\\mathcal{A}$ that solves $\\Pi$. It is known~\\cite{chang_pettie2019time_hierarchy_trees_rand_speedup,chang2020n1k_speedups} that randomness does not help for LCL problems on bounded-degree trees with round complexity $\\Omega(\\log n)$, so we assume that \n$\\mathcal{A}$ is deterministic.\n\nThe runtime of $\\mathcal{A}$ on a tree of size \n$O(w^{k})$ can be made to be at most $t = w\/10$ if $w$ is chosen to be sufficiently large. \n\nWe pick the labeling function $f$ as follows. \nWhenever we encounter a new bipolar tree $H$ to which we need to apply $f$, we simulate $\\mathcal{A}$ locally at the root $r^m$ of the middle rooted tree $T^m$ in $H$. Here we assume that the pumping lemma was applied before simulating $\\mathcal{A}$. Moreover, when we do the simulation, we assume that the number of vertices is $n = O(w^{k})$. \n\nTo make the simulation possible, we just locally generate arbitrary distinct identifiers in the radius-$t$ neighborhood $S$ of $r^m$. This is possible because $t = w\/10$ is so small that for any vertex $v$ in any tree $T$ constructed by recursively applying $Z_f$ for at most $k$ iterations, the radius-$(t+1)$ neighborhood of $v$ intersects at most one such $S$-set. Therefore, it is possible to complete the identifier assignment to the rest of the vertices in $T$ so that for any edge $\\{u,v\\}$, the set of identifiers seen by $u$ and $v$ are distinct in an execution of a $t$-round algorithm. \n\nThus, by the correctness of $\\mathcal{A}$, the labeling of all edges $\\{u,v\\}$ are configurations in $\\mathcal{E}$ and the labeling of all vertices $v$ are configurations in $\\mathcal{V}$.\nOur choice of $f$ implies that all trees in $W^\\ast$ admit correct labeling, and so $f$ is feasible. Combined with \\cref{lem:find_path_4}, we have proved that $\\mathsf{LOCAL}(n^{o(1)})$ $\\Rightarrow$ $\\ell$-full.\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\n\\section{Concluding Remarks}\\label{sec:concludingRemarksTrees}\n\nHere we collect some open problems we find interesting. \n\n\n\\begin{enumerate}\n \\item What is the relation of $\\mathsf{ffiid}, \\mathsf{MEASURE}, \\mathsf{fiid}$? In particular is perfect matching in the class $\\mathsf{ffiid}$ or $\\mathsf{MEASURE}$? \n \\item Are there LCL problems whose uniform complexity is nontrivial and slower than $\\operatorname{poly}\\log 1\/\\varepsilon$?\n \n \n \n \n \n \n \\item Suppose that $\\Pi_G$ is a homomorphism LCL, is it true that $\\Pi_G\\in\\mathsf{BOREL}$ if and only if $\\Pi_G\\in O(\\log^*(n))$?\n \\item Are there finite graphs of chromatic number bigger than $2\\Delta-2$ such that the corresponding homomorphism LCL is not in $\\mathsf{BOREL}$? See also Remark \\ref{r:hedet}.\n\\end{enumerate}\n\n\\subsection*{Acknowledgements} \nWe would like to thank Anton Bernshteyn, Endre Cs\u00f3ka, Mohsen Ghaffari, Jan Hladk\u00fd, Steve Jackson, Alexander Kechris, Edward Krohne, Oleg Pikhurko, Brandon Seward, Jukka Suomela, and Yufan Zheng for insightful discussions. \nWe especially thank Yufan Zheng for suggesting potentially hard instances of homomorphism LCL problems.\n\n\n\\bibliographystyle{alpha}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzkjno b/data_all_eng_slimpj/shuffled/split2/finalzzkjno new file mode 100644 index 0000000000000000000000000000000000000000..fbd0db0a1e08f7441019e48cacfe675f7f55e7ce --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzkjno @@ -0,0 +1,5 @@ +{"text":"\n\\section{Introduction}\n\\label{sec:intro}\n\\noindent The problem of wealth allocation is a key driver for past and present economic studies. Since the seminal contribution of \\cite{markowitz.1952}, both the industry and the researchers, have been highly involved in finding the optimal way to improve the performance of allocation strategies. During last decades, the Markowitz's approach has been widely criticised for the assumptions behind its simple mean--variance framework. In particular, the most critical points regard the homoskedasticity assumption and, in general, the Gaussian assumption for the joint returns distribution.\nRecently, researchers and practitioners found that the departure from the simple Markowitz's mean--variance approach, considering, for example, higher moments, can provide substantial improvements of the asset allocation strategies in terms of cumulative returns, see, e.g., \\cite{jondeau_rockinger.2003}, \\cite{lai_etal.2006}, \\cite{harvey_etal.2010} and \\cite{holly_etal.2011}. Furthermore, as for the conditional second moment, there have also been several attempts to model the conditional skewness and kurtosis of financial time series.\nIndeed, starting from the seminal contribution of \\cite{hansen.1994}, in the last two decades, the financial econometric literature has been focusing on modelling time--varying higher--order moments. One of the main issues is to figure out the correct way to introduce time--variation in skewness and kurtosis. Most important contributions in this sense are those of \\cite{hansen.1994}, \\cite{theodossiou_1998}, \\cite{harvey_siddique.1999}, \\cite{jondeau_rockinger.2003}, \\cite{deluca_etal.2006} and \\cite{jondeau_rockinger.2012}, who extend the model of \\cite{hansen.1994} and the class of ARCH models of \\cite{engle.1982} and \\cite{bollerslev.1986}. However, when dealing with higher--order conditional moments' dynamics in a multivariate context, the complexity increases principally for the high number of parameters needed and the major concern involves the overall tractability of the resulting model. As regards the co--skewness, main results are referred to \\cite{deluca_etal.2006}, \\cite{deluca_loperfido.2012} and \\cite{franceschini_loperfido.2010}, while for the co--kurtosis, up to our knowledge, there are not any contributions that impose a conditional dynamics on the fourth moment.\\newline\n\\indent The problem of dealing with higher order co--moments is strictly related to the general issue of describing the dependence structure affecting the multivariate financial time series. Since the seminal contribution of \\cite{sklar.1959}, copula models have been effectively used as an alternative flexible tool to the classical multivariate Gaussian models. Indeed, in the last years, copulas theory has played a central role for the financial econometric literature, as widely documented by \\cite{mcneil_etal.2005} and \\cite{joe.2014}. A comprehensive treatments of recent developments in this field can be found in the recent book of \\cite{durante_sempi.2015}. Furthermore, starting from the seminal contribution of \\cite{patton.2006a} and \\cite{jondeau_rockinger.2006b}, copula theory has been further extended to cope with dynamic evolving dependence structures. Time--varying copulas have been successfully employed in empirical works such as in \\cite{rodriguez.2007}, \\cite{delirasalvatierra_patton.2015} and \\cite{bernardi_catania.2015a}.\\newline\n\\indent An effective way to improve the models in sample fit as well as their forecast ability, is to include exogenous information. Nevertheless, most econometric models that aim to explain the joint dependence between economic factors, only consider endogenous information. By the way, it is natural to believe that, the dynamics affecting the relations between industries, and more generally between economies, are highly affected by exogenous, as well as, endogenous factors. It is worth noting that, the economic and financial literature focusing on univariate models, makes extensive use of exogenous information, see, e.g., \\cite{pagan.1971} and \\cite{baillie_1980} in the ARMAX context and \\cite{hwang_satchell.2005} for the GARCHX volatility model. Their empirical results confirm the superior ability of the models exploiting exogenous information, as compare with more simple alternatives. However, despite the documented usefulness of accounting for exogenous information in univariate time series models, in the financial econometric literature that focuses on the dynamic dependence modelling, there are only few examples of multivariate models dealing with exogenous information, see, e.g., \\cite{delirasalvatierra_patton.2015}. This lack is principally related to the difficulty of effectively including such information in such a way to retain model tractability, especially when the dimension of the observed variables is greater than 2. However, the contribution of including exogenous information in order to describe the dependence structure of economies or assets cannot be neglected. Let think about how the recent Global Financial Crises (GFC) of 2007--2008 influenced the dependence between economies. It is undoubtedly true that such a dramatic event drastically changed the choices of wealth allocation of economic agents. Indeed, capitals have moved away from the financial sector, and more generally from speculative instruments, such as options and futures, opening positions on less riskier assets such as short term government bonds or liquidity funds. In this reallocation of resources, it is obvious that, the behaviour of the sovereign interest rates dynamics, for examples, plays a central role in determining how much of the total wealth is reallocated, and hence how the dependence structure of the remaining assets is affected.\\newline\n\\indent In this paper, we propose a new Markov switching dynamic copula model being also able to deal with exogenous information in a fully multivariate setting using elliptical copulas. More precisely, the tractability induced by the copula methodology in dealing with marginal and joint separability is gathered with a new dynamic conditional dependence structure that relies on the presence of exogenous covariates and the inclusion of a latent Markovian structure, in order to enhance the model ability to account for time--varying higher moments. In this way, we are able to incorporate relevant information such as third and fourth conditional moments, as well as economic factors, in a portfolio optimisation framework. Throughout the paper, we will distinguish the marginal modelisation from the dependence structure linking individual asset into a multivariate framework. This distinction is available thanks to the use of copulas, and more precisely, to the \\cite{sklar.1959} Theorem and to the Inference Function for Margins (IFM) technique of \\cite{godambe_1960} and \\cite{mcleish_small.1988}.\nAs regards the univariate specification, we develop a new flexible parametric model accounting for time--varying conditional moments up to the fourth of the univariate distribution. In particular, we rely on the Generalised Autoregressive Score (GAS), also known as Dynamic Conditional Score (DCS), framework recently proposed by \\cite{creal_etal.2013} and \\cite{harvey.2013} in order to update the dynamics of time--varying parameters using the scaled score of the conditional distribution. This specification allows us to take account of a variety of well know stylised facts affecting the financial time series, like time varying moments, skewness and excess of kurtosis, as widely discussed by \\cite{mcneil_etal.2005}. Particularly interesting is the ability of the employed univariate model to approximate the unknown dynamics of the conditional density of returns by means of a filter based on the scaled score of the conditional distribution. This approach has been recently proved to be consistent with the goal of describing time variation for densities' parameters, see \\citep{koopman_etal.2015}.\n\\indent As stated before, the multivariate modelisation is carried on by using a copula function that links the dependence structure of each univariate series. This solution has been proved to provide highly flexible and reliable models for multivariate financial time series analysis, see e.g. \\cite{patton.2006b}, \\cite{joe.2014}, and the references quoted therein. Moreover, the seminal work of \\cite{patton.2006a}, also extend the original IFM two step estimator to the time--varying copula parameters framework. This fundamental result allows to effectively model the time--varying behaviour of linear as well as non--linear dependence patterns of financial returns. As discussed by \\cite{mcneil_etal.2005}, the flexibility in modelling linear and non--linear dependence patterns is one of the main stylised fact affecting multivariate financial time series. Indeed, as a direct consequence, econometric models dealing with multivariate financial time series, are required to account for this aspect, see, e.g., \\cite{delirasalvatierra_patton.2015}. For this reason, we propose a new specification for the evolution of the conditional correlation matrix of elliptical copulas relying on the famous Dynamic Conditional Correlation (DCC) specification of \\cite{engle.2002} and \\cite{tse_tsui.2002}, and the extension of \\cite{cappiello_etal.2006}. Moreover, as in \\cite{billio_caporin.2005} and, more recently, \\cite{bernardi_catania.2015a}, we also allow for the dynamic evolution of the conditional correlation matrix to depend upon the realisation of a first order Markovian process. More precisely, our model allows for different dynamic behaviours of the underlying dependence process in each specific state of the nature, as in \\cite{billio_caporin.2005} and \\cite{bernardi_catania.2015a}, but it additionally includes exogenous regressors in the dependence structure conditional to the latent state. The proposed model can effectively describe tranquil as well as turbulent periods affecting the financial markets, as for example, different phases of the economic or financial cycles. The opportunity to include exogenous covariates into the correlation dynamics is motivated by the observation that, in nowadays financial markets, economic agents usually base their decisions about wealth allocation conditionally on a huge amount of available information coming from multiple sources, which cannot be neglected. Summarising, our new specification is specifically tailored to model time variation in higher order conditional moments and possibly nonlinear time variation in the dependence structure of financial time series including the ability to account for exogenous information.\\newline\n\\indent The empirical part of the paper considers a panel of financial returns in a portfolio optimisation framework using higher order moments of the predictive conditional joint distribution. More precisely, in a similar way as \\cite{jondeau_rockinger.2006a}, we consider conditional portfolio moments into the expected utility of a rational investor by means of a Taylor expansion around the future wealth. Our empirical findings support our model since it outperform famous alternatives, like the DCC model. Specifically, the new specification improves the performance of allocation strategies based on the maximisation of the expected utility.\\newline\n\\indent The remaining of the paper is organised as follows. Section \\ref{sec:model} introduces the marginal and the joint dynamic model specifications and discusses their main features. Section \\ref{sec:asset_allocation} describes the asset allocation problem. Section \\ref{sec:empirical_study} deals with the empirical study and Section \\ref{sec:conclusion} concludes.\n\\section{The Model}\n\\label{sec:model}\n\\noindent In this paper, we propose to estimate the parameters governing the copula function and the associated dependence structure separately from those controlling for the dynamic evolution of the marginal distributions. This estimation technique, known as Inference Function for Marginals (IFM), has been proved to be feasible in large parametric spaces and asymptotically consistent for time varying copulas, see, e.g., \\cite{patton.2006a}. One of the main advantages of using the IFM estimating procedure concerns its ability to separate the univariate financial returns stylised facts from those regarding the multivariate distribution. This latter aspect, allows the researcher to focus on multivariate dependence structure only in a second moment, when the most appropriate model to describe all the empirical regularities of the observed univariate series has been selected and estimated. Whenever the marginal models are able to capture all the univariate empirical regularities of the marginal series, the only interestingly phenomenon to study is the remaining dependence structure. It follows that, the choice of an appropriate model that acts as a filter for the univariate series plays an important role for the subsequent multivariate analysis, see, e.g., \\cite{mcneil_etal.2005}.\\newline\n\\subsection{Marignal models}\n\\label{sec:marginal_model}\n\\noindent The financial econometrics literature offers several valid alternatives to model the conditional distribution of financial returns. In a fully parametric setting, following the same arguments presented by \\cite{cox.1981}, the first issue is to choose between the classes of observation driven and parameter driven models. The former can be well represented by the GARCH family of \\cite{engle.1982} and \\cite{bollerslev.1986}, while the latter is typically associated with the class of state space models; see for example \\cite{harvey_proietti.2005} and \\cite{durbin_koopman.2012}. However, the Generalised Autoregressive Score (GAS) framework, recently introduced by \\cite{creal_etal.2013} and \\cite{harvey.2013} is gaining lots of consideration by econometricians in many fields of time series analysis. Under the \\cite{cox.1981} classification, the GAS models can be considered as a class of observation driven models, with the usual consequences of having a closed form expression for the likelihood and ease of evaluation. However, as noted by \\cite{koopman_etal.2015}, the class of GAS models is similar to the class of state space models in the sense that both approaches make use of the information coming from the entire conditional distribution instead of using only that coming from the conditional expected value, as for example in the GARCH framework.\nIndeed, the key feature of GAS models is that the score of the conditional density is used as the forcing variable into the updating equation of a time--varying parameter. In the recent financial and econometrics literature, many reasons have been argued to support the adoption of score--based rules as a general updating mechanism of time--varying parameters, see, e.g., \\cite{harvey.2013} and \\cite{creal_etal.2013}, just to quote a few of them. Moreover, the flexibility of the GAS framework makes this class of models nested with a huge amount of famous econometrics models such as, for example, some of the ARCH--type models of \\cite{engle.1982} and \\cite{bollerslev.1986} for volatility modelling, and also the MEM, ACD and ACI models of \\cite{engle.2002b}, \\cite{engle_gallo.2006}, \\cite{engle_russell.1998} and \\cite{russell.1999}, respectively. Finally, one of the practical implications of using this framework in order to update the time--varying parameters, is that it avoids the problem of using a non--adequate forcing variable when the choice of it is not so obvious as for dynamic copula models; see, e.g., \\cite{bernardi_catania.2015a}.\\newline\n\\indent Formally, we assume that the $i$--th return at time $t$, for $t=1,2,\\dots,T$, is conditionally distributed according to the filtration $\\mathcal{F}_{i,t-1}=\\sigma\\left(y_{i,1},y_{i,2},\\dots,y_{i,t-1}\\right)$ as\n\\begin{equation}\ny_{i,t}\\sim\\mathcal{AST}\\left(y_{i,t};\\mu_{i,t},\\sigma_{i,t},\\gamma_{i,t},\\nu_{i,t}\\mid\\mathcal{F}_{t-1}\\right),\n\\end{equation}\nwhere $\\mathcal{AST}\\left(y_{i,t};\\cdot\\right)$ denotes the Asymmetric Student--t distribution with equal left and right tail decay of \\cite{zhu_galbraith.2010} with time varying location $\\mu_{i,t}\\in\\mathbb{R}$, scale $\\sigma_{i,t}\\in\\mathbb{R}^+$, shape $\\nu_{i,t}\\in\\left(4,\\infty\\right)$ and skewness $\\gamma_{i,t}\\in\\left(0,1\\right)$ parameters and density given by\n\\begin{equation}\nf_\\mathcal{AST}\\left(y_{i,t};\\mu_{i,t},\\sigma_{i,t},\\gamma_{i,t},\\nu_{i,t}\\right)=\\begin{cases}\n\\frac{1}{\\sigma_{i,t}}\\left[1+\\frac{1}{\\nu_{i,t}}\\left(\\frac{y_{i,t}-\\mu_{i,t}}{2\\gamma_{i,t}\\sigma_{i,t} K\\left(\\nu_{i,t}\\right)}\\right)\\right]^{-\\frac{\\left(\\nu_{i,t}+1\\right)}{2}},\\quad& y_{i,t}\\le\\mu_{i,t}\\nonumber\\\\\n\\frac{1}{\\sigma_{i,t}}\\left[1+\\frac{1}{\\nu_{i,t}}\\left(\\frac{y_{i,t}-\\mu_{i,t}}{2(1-\\gamma_{i,t})\\sigma_{i,t} K\\left(\\nu_{i,t}\\right)}\\right)\\right]^{-\\frac{\\left(\\nu_{i,t}+1\\right)}{2}},\\quad& y_{i,t}>\\mu_{i,t},\\nonumber\\\\\n\\end{cases}\n\\end{equation}\nwhere $K\\left(x\\right) = \\Gamma\\left(\\left(x+1\\right)\/2\\right)\/\\left[\\sqrt{\\pi x}\\Gamma\\left(x\/2\\right)\\right]$ and $\\Gamma\\left(\\cdot\\right)$ is the gamma function. We also define $\\mathbf{\\theta}_{i,t} = \\left(\\mu_{i,t},\\sigma_{i,t},\\gamma_{i,t},\\nu_{i,t}\\right)^\\prime\\in\\mathbb{R}\\times\\mathbb{R}^+\\times\\left(0,1\\right)\\times\\left(4,\\infty\\right)$ to be a vector containing all the time--varying parameters. Furthermore, let $h:\\mathbb{R}^4\\to\\mathbb{R}\\times\\mathbb{R}^+\\times\\left(0,1\\right)\\times\\left(4,\\infty\\right)$ be a continuous, $\\mathcal{F}_{t-1}$ measurable, twice differentiable, vector--valued mapping function, such that $\\mathbf{\\theta}_t = h\\left(\\tilde{\\mathbf{\\theta}}_{i,t}\\right)$, where $\\tilde\\mathbf{\\theta}_{i,t}\\in\\mathbb{R}^4$ is a proper reparametrisation of $\\mathbf{\\theta}_{i,t}$, for all $t=1,2,\\dots,T$. In our context, a convenient choice for $h\\left(\\cdot\\right)$ is\n\\begin{equation}\nh\\left(\\tilde{\\mathbf{\\theta}}_{i,t}\\right) = \\begin{cases}\n\\mu_{i,t}=\\mu_{i,t}\\\\\n\\sigma_{i,t}=\\exp\\left(\\tilde\\sigma_{i,t}\\right)\\\\\n\\gamma_{i,t}=\\left[1 + \\exp\\left(-\\tilde\\gamma_{i,t}\\right)\\right]^{-1}\\\\\n\\nu_{i,t}=4+\\exp\\left(\\tilde\\nu_{i,t}\\right).\n\\end{cases}\n\\end{equation}\nWe let the time--varying reparameterised vector $\\tilde\\mathbf{\\theta}_{i,t}=\\left(\\mu_{i,t},\\tilde\\sigma_{i,t},\\tilde\\gamma_{i,t},\\tilde\\nu_{i,t}\\right)^\\prime$ to be updated using the score of the conditional distribution of $y_{i,t}$, exploiting the GAS dynamic where the Fisher information matrix of the conditional distribution is used as scaling matrix, as suggested by \\cite{creal_etal.2013} and \\cite{harvey.2013}. Specifically, we consider the following updating mechanism for the parameter dynnamics\n\\begin{align}\n \\begin{pmatrix}\n \\mu_{i,t+1}\\\\\n \\tilde\\sigma_{i,t+1}\\\\\n \\tilde\\gamma_{i,t+1}\\\\\n \\tilde\\nu_{i,t+1}\n \\end{pmatrix} =\n \\begin{pmatrix}\n \\omega_{\\mu_i}\\\\\n \\omega_{\\sigma_i}\\\\\n \\omega_{\\gamma_i}\\\\\n \\omega_{\\nu_i}\n \\end{pmatrix} +\n \\begin{pmatrix}\n \\alpha_{\\mu_i} & 0 & 0 & 0\\\\\n 0 & \\alpha_{\\sigma_i} & 0 & 0\\\\\n 0 & 0 & \\alpha_{\\gamma_i} & 0 \\\\\n 0 & 0 & 0 & \\alpha_{\\nu_i}\n \\end{pmatrix} \\tilde{s}_{i,t} +\n \\begin{pmatrix}\n \\beta_{\\mu_i} & 0 & 0 & 0\\\\\n 0 & \\beta_{\\sigma_i} & 0 & 0\\\\\n 0 & 0 & \\beta_{\\gamma_i} & 0 \\\\\n 0 & 0 & 0 & \\beta_{\\nu_i}\n \\end{pmatrix}\n \\begin{pmatrix}\n \\mu_{i,t}\\\\\n \\tilde\\sigma_{i,t}\\\\\n \\tilde\\gamma_{i,t}\\\\\n \\tilde\\nu_{i,t}\n \\end{pmatrix},\n \\label{eq:gas_empDyn}\n\\end{align}\nwhere $\\tilde s_t = \\mathcal{J}\\left(\\tilde\\mathbf{\\theta}_{i,t}\\right)^{-1}\\mathcal{I}\\left(\\mathbf{\\theta}_{i,t}\\right)^{-1}\\nabla\\left(y_{i,t};\\mathbf{\\theta}_{i,t}\\right)$ is the scaled score of the reparameterised conditional distribution of $y_{i,t}$. The quantity $\\tilde s_t$ depends on the product of three matrices defined as\n\\begin{align}\n\\mathcal{J}\\left(\\tilde\\mathbf{\\theta}_{i,t}\\right) &= \\frac{\\partial\\mathrm{h}\\left(\\tilde{\\boldsymbol{\\theta}}_{i,t}\\right)}{\\partial\\tilde\\mathbf{\\theta}_{i,t}}\\\\\n\\mathcal{I}\\left(\\mathbf{\\theta}_{i,t}\\right) & = \\mathbb{E}\\left[\\nabla\\left(y_{i,t};\\mathbf{\\theta}_{i,t}\\right)\\times\\nabla\\left(y_{i,t};\\mathbf{\\theta}_{i,t}\\right)^\\prime\\mid\\mathcal{F}_{t-1}\\right]\\\\\n\\nabla\\left(y_{i,t};\\mathbf{\\theta}_{i,t}\\right) &= \\frac{\\partial\\ln f_\\mathcal{AST}\\left(y_{i,t};\\mathbf{\\theta}_{i,t}\\right)}{\\partial{\\mathbf{\\theta}_{i,t}}},\n\\end{align}\nwhere $\\mathcal{J}\\left(\\tilde\\mathbf{\\theta}_{i,t}\\right)$ is the Jacobian of the mapping function $h\\left(\\cdot\\right)$, i.e.\n\\begin{equation}\n \\mathcal{J}\\left(\\tilde\\mathbf{\\theta}_{i,t}\\right) = \\begin{pmatrix}\n 1 & 0 & 0 & 0 \\\\\n 0 & \\exp\\{\\tilde\\sigma_{i,t}\\} & 0 & 0 \\\\\n 0 & 0 & -\\exp\\{-\\tilde\\gamma_{i,t}\\}\\left(1+\\exp\\{-\\tilde\\gamma_{i,t}\\}\\right)^{-2} & 0 \\\\\n 0 & 0 & 0 & \\exp\\{\\tilde\\nu_{i,t}\\},\n \\end{pmatrix},\n\\end{equation}\nwhere $\\mathcal{I}\\left(\\mathbf{\\theta}_{i,t}\\right)$ and $\\nabla\\left(y_{i,t};\\mathbf{\\theta}_{i,t}\\right)$ are the Fisher information matrix and the score of the conditional AST distribution, respectively. For the parameterisation here considered, the Fisher information matrix of the AST distribution is reported in \\cite{zhu_galbraith.2010}, while the score is reported in \\ref{sec:ScoreAndFisher}.\\newline\n\\indent It is worth noting that, in equation \\eqref{eq:gas_empDyn}, a diagonal structure for the matrix of coefficient is considered. Obviously, this is a convenient choice in order to reduce the number of model parameters, but other solutions are possible. In order to guarantee the stationarity of the process defined in equation \\eqref{eq:gas_empDyn} we only need to impose the constraint that all the autoregressive coefficients $\\beta_k$, $k\\in\\left(\\mu,\\sigma,\\gamma,\\nu\\right)$ would be in modulus less then one, since, under a correct model specification, the sequence of scores $\\{\\tilde s_t, t>0\\}$ is a martingale difference.\n\\subsection{Multivariate model}\n\\label{sec:multivariate_model}\n\\noindent Dynamic time--varying dependence structures of financial assets have been deeply investigated in the recent developments of financial econometrics. One of the most successful and powerful model has been the Dynamic Conditional Correlation (DCC) model of \\cite{engle.2002} and \\cite{tse_tsui.2002} which allows for time variation in the correlation of assets returns. In the last few years many models were also been proposed in order to improve the flexibility of the original DCC, see, e.g., \\cite{bauwens_etal.2006} for a complete surveys of these recent developments. The success of the DCC model is due to its two stage estimation procedure which mainly relies on the multivariate Gaussian assumption. However, when the Gaussian assumption is violated, the separability of the log likelihood does not hold anymore and the only viable alternative in order to remain with a feasible model is to adopt a Quasi--MLE estimation approach, see, e.g., \\cite{engle_sheppard.2001}. Despite its enormous popularity, only few works have been focused on the inclusion of exogenous information in a DCC framework in order to let the covariance co--movements to depend on covariates. Among those, \\cite{vargas.2008} firstly introduced exogenous covariates into the Asymmetric Generalised DCC (AGDCC) model of \\cite{cappiello_etal.2006}, while \\cite{chou_liao.2008} follow a similar approach on a simpler DCC model. However, both the empirical applications of the aforementioned papers, only consider a bivariate setting under the Gaussian assumption for the innovation terms. Besides the tractability of the resulting model, one of the main difficulties of including exogenous covariates into a DCC model, is to guarantee the positiveness of the variance--covariance matrix at each point in time. A common solution relies on the Constrained Maximum Likelihood (CML) or Penalised Maximum Likelihood (PML) estimation approaches. In the context of bivariate time--varying copulas, only \\cite{delirasalvatierra_patton.2015} consider the inclusion of exogenous information to model the correlation matrix dynamic. However, \\cite{delirasalvatierra_patton.2015} do not consider a DCC--type dynamic evolution of the copula parameters, but, they instead rely on the GAS framework of \\cite{creal_etal.2013} and \\cite{harvey.2013}. Moreover, their model cannot be used in a fully multivariate setting when considering elliptical copulas such as the Gaussian and the Student--t. In what follows, we introduce our Student--t dynamic Markov--Switching (MS) copula model where the copula dependence matrix evolves according to a MS--DCC model which also depends on exogenous covariates.\\newline\n\\indent Let $\\mathbf{u}_t=\\left(u_{1,t},u_{2,t},\\dots,u_{N,t}\\right)^\\prime$ with $u_{j,t}=F_{\\mathcal{AST}}\\left(y_{i,t};\\mathbf{\\theta}_{i,t}\\right)$ be the Probability Integral Transformation (PIT) of $y_{i,t}$ according to the conditional AST distribution $F_{\\mathcal{AST}}\\left(\\cdot\\right)$ with parameters $\\mathbf{\\theta}_{i,t}=\\left(\\mu_{i,t},\\sigma_{i,t},\\gamma_{i,t},\\nu_{i,t}\\right)^\\prime$, and assume that\n\\begin{equation}\n\\mathbf{u}_t\\mid S_t=s\\sim\\mathcal{C}_T\\left(\\mathbf{u}_t;\\mathbf{R}_t^s,\\nu_\\mathsf{c}^s\\right),\n\\end{equation}\nfor $t=1,2,\\dots,T$, where $\\mathbf{R}_t^s$ is the time--varying, regime dependent correlation matrix at time $t$ of the pseudo--observation $\\mathbf{u}_t$, and $\\nu_\\mathsf{c}^s$ is the degree of freedom parameter that is also subject to the realisation of the first order Markov chain. According to the DCC dynamics the conditional correlation matrix $\\mathbf{R}_t^s$ can be decomposed in the following way\n\\begin{align}\n\\mathbf{R}_t^s={\\mathbf{D}_t^s}^{-1}\\mathbf{C}_t^s{\\mathbf{D}_t^s}^{-1},\n\\end{align}\nwhere $\\mathbf{D}_t^{s}$ is a diagonal matrix containing the square root of the diagonal elements of $\\mathbf{C}_t^s$. We further assume that the latent states $s=1,2,\\dots,L$ are driven by a Markov process $S_t$, for $t=1,2,\\dots,T$ defined on the discrete space $\\Omega=\\left\\{1,2,\\dots,L\\right\\}$ with transition probability matrix $\\mathbf{Q}=\\left\\{q_{l,k}\\right\\}$, where $q_{l,k}=\\mathbb{P}\\left(S_t=k\\mid S_{t-1}=l\\right)$, $\\forall l,k\\in\\Omega$ is the probability that state $k$ is visited at time $t$ given that at time $t-1$ the chain was in state $l$, and initial probabilities vector $\\mbox{\\boldmath $\\delta$}=\\left(\\delta_1,\\delta_2,\\dots,\\delta_L\\right)^\\prime$, $\\delta_l=\\mathbb{P}\\left(S_1=l\\right)$, i.e., the probability of being in state $l=\\left\\{1,2,\\dots,L\\right\\}$ at time 1. For an up to date review of HMMs, see, e.g., \\cite{cappe_etal.2005}, \\cite{zucchini_macdonald.2009} and Dymarski \\citeyearpar{dymarski.2011}. The following dynamics is imposed on the state dependent variance--covariance matrix\n\\begin{align}\n\\mathbf{C}_{t+1}^s&= \\left(\\bar{\\mathbf{C}}-\\mathbf{A}^s\\bar{\\mathbf{C}}{\\mathbf{A}^{s}}^\\prime-\\mathbf{B}^s\\bar{\\mathbf{C}}{\\mathbf{B}^{s}}^\\prime-\\mathbf{K}{\\mbox{\\boldmath $\\xi$}^{s}}^\\prime\\bar{\\mathbf{X}}\\right)\n+\\mathbf{A}^s\\bXi_t{\\mathbf{A}^{s}}^\\prime+\\mathbf{B}^s\\mathbf{C}_t^s{\\mathbf{B}^{s}}^\\prime+{\\mathbf{K}\\mbox{\\boldmath $\\xi$}^{s}}^\\prime\\mathbf{X}_t,\n\\label{eq:corr_dynamic}\n\\end{align}\nfor $s=1,2,\\dots,L$, where $\\mathbf{A}^s$ and $\\mathbf{B}^s$ are diagonal matrices with elements $\\mathbf{A}^s=\\left\\{\\sqrt{\\alpha_i^s}\\right\\}$, $\\mathbf{B}^s=\\left\\{\\sqrt{\\beta_i^s}\\right\\}$, for $i=1,2,\\dots,N$, $s=1,2,\\dots,L$, $\\mathbf{K}$ is a $N\\times N$ matrix of ones, $\\mbox{\\boldmath $\\xi$}^s$ is a $p\\times 1$ vector containing the coefficients associated to the exogenous variables $\\mathbf{X}_t$. Here, $\\bar\\mathbf{C}$ denotes the empirical variance--covariance matrix of the pseudo--observations $\\mathbf{u}_t$ and $\\bar\\mathbf{X}$ is a $p\\times 1$ vector containing the empirical mean of the regressors, which are introduced in equation \\eqref{eq:corr_dynamic} in order to enforce the mean--reversion of the process. Moreover, in order to guarantee the stationarity of the process defined in equation \\eqref{eq:corr_dynamic}, we impose $\\alpha_i^s + \\beta_i^s<1$ for all $i=1,2,\\dots,N$ and $s=1,2,\\dots,L$. Moreover, we need to ensure that $\\mathbf{C}_{t}^s$ is positive defined matrix for each $t=1,2,\\dots,T$ and across all the possible states of nature $s=1,2,\\dots,L$.\nThe only unspecified quantity in equation \\eqref{eq:corr_dynamic}, is the forcing variable $\\bXi_t$ which is set to the moving average variance--covariance matrix of length $m$\n\\begin{equation}\n\\bXi_t = m^{-1}\\sum_{j=0}^m \\mathbf{u}_{t-m+j}\\mathbf{u}_{t-m+j}^\\prime - \\bar\\mathbf{u}_{m,j}\\bar\\mathbf{u}_{m,j}^\\prime,\n\\end{equation}\nwhere $\\bar\\mathbf{u}_{m,j} = m^{-1}\\sum_{j=0}^m \\mathbf{u}_{t-m+j}$ represents the vector containing the simple mean of the pseudo observations across the period $\\left\\{t-m;t\\right\\}$. The proposed model is similar those of \\cite{vargas.2008} and \\cite{chou_liao.2008} but differs for the chosen forcing variable $\\bXi_t$ which updates the conditional variance--covariance matrix. The proposed updating scheme makes use of a\nrolling window variance--covariance matrix as forcing variable by which we can get a smooth time evolution of the correlation.\nAs regards this aspect, the proposed correlation dynamics generalises the approach of \\cite{patton.2006b} and \\cite{jondeau_rockinger.2006b} in a fully multivariate setting.\\newline\n\\indent More interesting, the model here proposed, allows for the dynamic of the variance--covariance matrix to be also subject to the realisation of a first order Markov process. Using a different approach, \\cite{bernardi_catania.2015a} show that including this regime dependent behaviour into the variance--covariance dynamics, really helps in describing the time--varying dependence of equity indexes. Indeed, ad documented, for example, by \\cite{pelletier.2006}, \\cite{guidolin_timmermann.2008}, \\cite{ang_bekaert.2002}, \\cite{bernardi_catania.2015a}, there is strong evidence of the presence of different regimes affecting the dependence structure of equity assets. Obviously, the model proposed in equation \\eqref{eq:corr_dynamic} can be easily generalised to account for the leverage effect in a similar way as for the AGDCC of \\cite{cappiello_etal.2006}\n\\begin{align}\n\\mathbf{C}_{t+1}^s =&\\left(\\bar\\mathbf{C}-\\mathbf{A}^s\\bar\\mathbf{C}{\\mathbf{A}^{s}}^\\prime-\\mathbf{B}^s\\bar\\mathbf{C}{\\mathbf{B}^{s}}^\\prime-\\mbox{\\boldmath $\\Gamma$}^s\\mathrm{\\bar\\mathbf{N}}^s{\\mbox{\\boldmath $\\Gamma$}^{s}}^\\prime-\\mathbf{K}{\\mbox{\\boldmath $\\xi$}^{s}}^\\prime\\bar\\mathbf{X} \\right)\\nonumber\\\\\n&\\qquad\\qquad\\qquad+\\mathbf{A}^s\\bXi_t{\\mathbf{A}^{s}}^\\prime+\\mathbf{B}^s\\mathbf{C}_t^s{\\mathbf{B}^{s}}^\\prime+\\mbox{\\boldmath $\\Gamma$}^s\\boldsymbol\\eta_t^s{\\boldsymbol\\eta_t^s}^\\prime{\\mbox{\\boldmath $\\Gamma$}^{s}}^\\prime+\\mathbf{K}{\\mbox{\\boldmath $\\xi$}^{s}}^\\prime\\mathbf{X}_t,\n\\label{eq:Gdynamic}\n\\end{align}\nwhere $\\mbox{\\boldmath $\\Gamma$}^s=\\left\\{\\sqrt{\\gamma_i^s}\\right\\}$, $i=1,2,\\dots,N$ is a diagonal matrix, and $\\boldsymbol\\eta_t^s$ is a vector whose $i$--th component is given by $\\bbone\\left(x_{i,t}^s<0\\right)$ where $x_{i,t}^s=\\mathcal{T}_{\\nu_\\mathsf{c}^s}^{-1}\\left(u_{it}\\right)$ for $i=1,2,\\dots,N$, and $\\mathcal{T}_{\\nu_\\mathsf{c}^s}^{-1}\\left(\\cdot\\right)$ denotes the inverse cumulative density function (cdf) of a standardised Student--t distribution with $\\nu_\\mathsf{c}^s$ degree of freedom, and $\\bbone\\left(\\cdot\\right)$ represents the indicator function.\nIn equation \\eqref{eq:Gdynamic}, $\\mathrm{\\bar\\mathbf{N}}^s$ is the unconditional empirical average of the matrices $\\boldsymbol\\eta_t^s{\\boldsymbol\\eta_t^s}^\\prime$, $t=1,2,\\dots,T$. Several constraints must be imposed in order to avoid explosive patterns in the conditional variance--covariance matrix dynamic in equation \\eqref{eq:Gdynamic}, see, e.g., \\cite{cappiello_etal.2006}. Since in our empirical investigation we do not find any substantial improvement in favour of model \\eqref{eq:Gdynamic} with respect to model \\eqref{eq:corr_dynamic}, throughout the paper we will continue to refer to the specification without leverage effect. Moreover, since the number of parameters in the model in described in equation \\eqref{eq:corr_dynamic} grows linearly with respect to the number of asset $N$, sometimes it would be convenient to set $\\alpha_1^s = \\alpha_2^s = \\dots = \\alpha_N^s=\\alpha^s$ and $\\beta_1^s = \\beta_2^s = \\dots = \\beta_N^s=\\beta^s$. We name this constrained specification `` simple'', while we refer to the aforementioned general case with the name `` generalised''.\\newline\n\\indent To conclude the model specification, we report the multivariate density function of the random vector $\\mathbf{y}_t=\\left(y_{1,t},y_{2,t},\\dots,y_{N,t}\\right)^\\prime$ which is given by\n\\begin{align}\n\\mathbf{y}_{t+1}\\mid\\mathbf{y}_{t}\\sim h\\left(\\mathbf{y}_{t+1};\\mathbf{\\theta}_{t+1},\\mathbf{X}_{t},\\mbox{\\boldmath $\\Delta$}_{t+1}^s, s=1,\\dots,L\\right),\n\\label{eq:multivariate_density}\n\\end{align}\nwhere $\\mathbf{\\theta}_t=\\left(\\mathbf{\\theta}_{1,t}^\\prime,\\mathbf{\\theta}_{2,t}^\\prime,\\dots,\\mathbf{\\theta}_{N,t}^\\prime\\right)^\\prime$, and\n\\begin{align}\nh\\left(\\mathbf{y}_{y+1};\\cdot\\right)=&\\prod_{i=1}^N f_{\\mathcal{AST}}\\left(y_{i,t+1};\\mathbf{\\theta}_{i,t+1}\\right)\n\\sum_{s=1}^L\\pi_{t+1\\vert t}^{(s)} c_T\\left(F_{\\mathcal{AST}}\\left(\\mathbf{y}_{t+1};\\mathbf{\\theta}_{i,t+1}\\right),\\mbox{\\boldmath $\\Delta$}_{t+1}^s,\\mathbf{X}_t\\right)\\nonumber,\n\\end{align}\nwith $\\mbox{\\boldmath $\\Delta$}_{t+1}^s=\\left(\\mathbf{R}_{t+1}^s,\\nu_{\\mathsf{c}}^s\\right)$ and the mixing weight $\\pi_{t+1\\vert t}^{(s)}$ are\n\\begin{equation}\n\\pi_{t+1\\vert T}^{(s)}=\\sum_{m=1}^L {q}_{m,s}\\mathbb{P}\\left(S_{t+1}=m\\mid\\mathbf{y}_{1:t},\\mathbf{X}_{1:t}\\right),\\nonumber\n\\end{equation}\nfor $s=1,2,\\dots,L$, and $\\mathbb{P}\\left(S_{t+1}=m\\mid\\mathbf{r}_{1:t},\\mathbf{X}_{1:t}\\right)$ can be evaluated using the well known FFBS algorithm detailed in \\cite{fruhwirth_schnatter.2006}. Here, the $i$--th marginal density and probability functions are represented by $f_\\mathcal{AST}\\left(\\cdot\\right)$ and $F_\\mathcal{AST}\\left(\\cdot\\right)$ while the copula probability density is denoted by $c_T\\left(\\cdot\\right)$.\\newline\n\\indent Concerning the estimation of the proposed model, unfortunately, given the presence of the exogenous covariates, we cannot rely on the usual results for DCC models such as, for examples, those of \\cite{engle.2002} and \\cite{cappiello_etal.2006}, while, the Constrained ML (CML) approach used by \\cite{vargas.2008} and \\cite{chou_liao.2008} in a similar context, requires the evaluation of highly nonlinear constraint during the maximisation procedure. We follow a different approach which consists on penalising the log--likelihood function (PML) in the second step of the IFM procedure, whenever the variance--covariance matrix is not positive defined. However, this solution is not optimal since it introduces strongly nonlinearities in the likelihood shape, and it may result in possible local optima solutions of the maximisation problem. Hence, good starting values are required.\nIn the second step of the IFM procedure, we made use of the Expectation--Maximisation algorithm of \\cite{dempster_etal.1977} in order to deal with the markovian structure characterising the latent states. Further details about the employed estimation technique can be found in \\cite{bernardi_catania.2015a}.\n\\section{Asset Allocation Strategy}\n\\label{sec:asset_allocation}\n\\noindent In this section, we tailor the `` distribution timing\\qmcsp approach to the portfolio optimisation problem of \\cite{jondeau_rockinger.2012} to the dynamic MS copula approach introduced in the previous Section. Specifically, we consider a rational investor having full information about the whole distribution of his\/her future wealth and builds a dynamic asset allocation strategy by maximising his\/her expected utility. It is worth noting that, higher moments play a fundamental role during the utility maximisation process.\nFormally, let $y_{p,t+1} = \\sum_{i=1}^N\\lambda_{i,t\\vert t+1} y_{i,t+1}$ be the future portfolio return, where $\\lambda_{i,t\\vert t+1}$ denotes the portion of today wealth that the investor is willing to allocate in the $i$--th risky asset, with $\\sum_{i=1}^N \\lambda_{i,t\\vert t+1}=1$, then $W_{t+1} = 1 + y_{p,t+1}$ represents the future portfolio wealth at time $t+1$\\footnote{We will always consider an initial wealth of 1\\$. This choice will not affect the portfolio decision given our assumption on the investor's utility function.}. The vectors, $\\mathbf{y}_{t+1} = \\left(y_{1,t+1},y_{2,t+1},\\dots,y_{N,t+1}\\right)^\\prime$ and $\\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}=\\left(\\lambda_{1,t\\vert t+1},\\lambda_{2,t\\vert t+1},\\dots,\\lambda_{N,t\\vert t+1}\\right)^\\prime$ contain all the assets in which the investor can open a speculative position at time $t$, which is subsequently held for the period $\\left[t,t+1\\right)$ and the portfolio weights, respectively. The portfolio optimisation problem can then be written as follows\n\\begin{equation}\n\\arg\\max_{\\left\\{\\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}\\right\\}}\\mathbb{E}_t\\left[\\mathcal{U}\\left(1 + \\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}^\\prime \\mathbf{y}_{t+1}\\right)\\right],\\nonumber\n\\label{eq:maximisation}\n\\end{equation}\nwhere $\\mathcal{U}\\left(W_{t+1}\\right)$ denotes the investor utility function and the expectation must be taken with respect to the future portfolio return distribution. Clearly, this problem does not have a closed--form solution given that $\\mathbf{y}_{t+1}$ is conditionally distributed according to equation \\eqref{eq:multivariate_density}. Moreover, when $N$ is large, numerical techniques such as, for example, Monte Carlo integration and quadrature rules become infeasible. Nevertheless, the maximisation problem can be easily solved by taking a Taylor expansion of the utility function around the wealth at time $t$, $W_t$, i.e., by considering\n\\begin{equation}\n\\mathcal{U}\\left(W_{t+1}\\right) = \\sum_{k=0}^\\infty\\frac{\\mathcal{U}^{\\left(k\\right)}\\left(W_t\\right)}{k!}\\left(W_{t+1} - W_t\\right)^k,\n\\label{eq:taylor_expansion}\n\\end{equation}\nwhere $\\mathcal{U}^{\\left(k\\right)}$ denotes the $k$--th derivative of the utility function with respect to its argument. Taking the expectation at time $t$ of both sides of equation \\eqref{eq:taylor_expansion} we get\n\\begin{equation}\n\\mathbb{E}_t\\left[\\mathcal{U}\\left(W_{t+1}\\right)\\right] =\\sum_{k=0}^\\infty\\frac{\\mathcal{U}^{\\left(k\\right)}\\left(W_t\\right)}{k!}m_{p,t+1}^{\\left(k\\right)},\n\\label{eq:taylor_expansion_portfolio}\n\\end{equation}\nwhere $m_{p,t+1}^{\\left(k\\right)}=\\mathbb{E}\\left[r_{p,t+1}^k\\right]=\\mathbb{E}\\left[\\left(W_{t+1}-W_t\\right)^k\\right]$ denotes the $k$--th non central moment of the portfolio return distribution at time $t+1$. If we consider a fourth--order Taylor expansion, i.e., we truncate equation \\eqref{eq:taylor_expansion_portfolio} at $\\bar{k}=4$, the first four central moments of the portfolio distribution are immediately available as a function of the corresponding non--central moments and the portfolio weights. More precisely, the centred moments of the portfolio returns, at time $t+1$, are\n\\begin{align}\n\\mu_{p,t+1} &= \\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}^\\prime M_{1,t+1}\\nonumber\\\\\n\\sigma_{p,t+1}^2 &= \\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}^\\prime M_{2,t+1}\\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}\\nonumber\\\\\ns_{p,t+1}^3 &= \\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}^\\prime M_{3,t+1}\\left(\\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}\\otimes \\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}\\right)\\nonumber\\\\\nk_{p,t+1}^4 &= \\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}^\\prime M_{4,t+1}\\left(\\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}\\otimes \\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}\\otimes \\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}\\right),\\nonumber\n\\end{align}\nwhere $\\otimes$ stands for the Kronecker product, and the corresponding non--centred moments required by equation \\eqref{eq:taylor_expansion_portfolio} can be obtained as\n\\begin{align}\nm_{p,t+1}^{\\left(1\\right)} &= \\mu_{p,t+1} \\nonumber\\\\\nm_{p,t+1}^{\\left(2\\right)} &= \\sigma_{p,t+1}^2 + \\mu_{p,t+1}^2 \\nonumber\\\\\nm_{p,t+1}^{\\left(3\\right)} &= s_{p,t+1}^3 + 3\\sigma_{p,t+1}^2\\mu_{p,t+1} + \\mu_{p,t+1}^3\\nonumber\\\\\nm_{p,t+1}^{\\left(4\\right)} &= k_{p,t+1}^4 + 4s_{p,t+1}^3\\mu_{p,t+1} + 6\\sigma_{p,t+1}^2\\mu_{p,t+1}^2 + \\mu_{p,t+1}^4.\\nonumber\n\\end{align}\nNote that $M_{1,t+1}$, $M_{2,t+1}$, $M_{3,t+1}$, $M_{4,t+1}$ are the mean, the variance--covariance matrix, the co--skewness and co--kurtosis of the conditional joint density function in equation \\eqref{eq:multivariate_density}, respectively. Here, the original $\\left(N,N,N\\right)$ and $\\left(N,N,N,N\\right)$ matrices of co--skewness and co--kurtosis were transformed into the $\\left(N,N^2\\right)$ and $\\left(N,N^3\\right)$ matrices, respectively, as suggested by \\cite{deathayde_flores.2004} by slicing each $\\left(N,N\\right)$ layer and pasting them, in the same order, sideways.\\newline\n\\indent Having defined the general portfolio optimisation framework, the only quantity that remains unspecified is the specific utility function the rational investor maximises. We assume that our agent makes his\/her investment choices according to a Constant Relative Risk Aversion (CRRA) utility function defined as\n\\begin{equation}\n\\mathcal{U}_{t+1}\\left(W_{t+1}\\right) = \\begin{cases} \\frac{W_{t+1}^{1-\\upsilon}}{1-\\upsilon}, & \\mbox{if } \\upsilon>1 \\\\ \\log\\left(W_{t+1}\\right), & \\mbox{if } \\upsilon=1, \\end{cases}\n\\label{eq:utility}\n\\end{equation}\nwhere $\\upsilon$ is the coefficient of relative risk aversion, which is assumed to be constant with respect to the wealth $W_{t+1}$. For CRRA utility function, equation \\eqref{eq:taylor_expansion_portfolio} truncated at the fourth order becomes\n\\begin{align}\n\\mathbb{E}_t\\left[\\mathcal{U}\\left(W_{t+1}\\right)\\right]&\\approx\\frac{1}{1-\\upsilon} + m_{p,t+1}^{\\left(1\\right)} - \\frac{\\upsilon}{2} m_{p,t+1}^{\\left(2\\right)}\n+\\frac{\\upsilon\\left(\\upsilon + 1\\right)}{6} m_{p,t+1}^{\\left(3\\right)} -\n\\frac{\\upsilon\\left(\\upsilon + 1\\right)\\left(\\upsilon + 2\\right)}{24} m_{p,t+1}^{\\left(4\\right)},\n\\label{eq:utility_final}\n\\end{align}\nwhere the approximation order is $O\\left(5\\right)$. The closed--form expression for the expected utility makes the maximisation of equation \\eqref{eq:utility_final} with respect to the portfolio weights, straightforward. However, unfortunately, there is no closed--form solution for the moments of order two to four of the multivariate joint density function and we need to resort to numerical procedures. \\ref{sec:appendix_Moments} describes the method used to approximate the moments of the predictive joint density.\n\\section{Empirical study}\n\\label{sec:empirical_study}\n\\noindent In this section we apply the model and the portfolio optimisation methodology introduced in previous sections to deliver a set of optimal portfolio weights. To this end, we consider the same data set of \\cite{jondeau_rockinger.2006a} which consists of five of the major international equity indexes, namely the Standard \\& Poors 500 (SPX), the Nikkei 225 (N225), the FTSE 100 (FTSE), the DAX30 (GDAXI), and the CAC40 (FCHI). At each time $t$ over the out of sample period, the rational investor holding the portfolio is expected to choose the optimal allocation of the wealth by maximising his\/her expected utility using the argument presented in Section \\ref{sec:asset_allocation}. The illustration of the empirical results firstly focus on the marginal distribution, then deals with the multivariate density estimation and forecast and concludes with the presentation of the optimal portfolio allocations.\n\\subsection{Data}\n\\label{sec:data}\n\\noindent The data set consists of 1449 weekly returns spanning the period from the 8--th January, 1988, to the 9--th October, 2015. Starting from this series, the last 449 observations spanning the period from the 9--th March, 2007 to the end of the sample, are used to perform the out of sample analysis and to backtest the portfolio strategy, while the fist 1000 observation are used to fit the model and assess its in--sample performance. Table \\ref{tab:Index_data_summary_stat} reports descriptive statistics of the indexes returns over the in--sample and out--of--sample periods. As expected, we find strong evidence of departure from normality, since all the considered series appear to be leptokurtic and skewed. Moreover, the \\cite{jarque_bera.1980} statistic test strongly rejects the null hypothesis of normality for all the considered series at the 1\\% confidence level. It is worth noting that, the departure from normality, appears to be stronger during the out--of--sample period. As discussed for example by \\cite{shiller_2012}, this empirical evidence can be considered as an effect of the recent GFC of 2007--2008 that affected the overall economy. Table \\ref{tab:Dependence} reports the linear correlation and the Kendall $\\tau$ unconditional estimates. All the coefficients seem to be substantially different between the two subperiods, suggesting a time varying behaviour of the dependence structure that characterise our series. To further investigate this aspect, we perform the LMC test for constant versus time--varying conditional correlation of \\cite{tse.2000} on the returns belonging to the whole period. In our case the LMC test is distributed according to a $\\chi_{10}^2$ with a critical value of about 23.21 at the 1\\% confidence level, while the test statistic is about 997 which strongly adverses the null hypothesis of time invariant correlation. Moreover, the lower triangular parts of Table \\ref{tab:Dependence} also gives an insights about the changes in the tail of the joint probability distribution. These findings strongly support for the Student--t copula DCC model with exogenous regressors introduced in the previous sections \\ref{sec:marginal_model}--\\ref{sec:multivariate_model}, being able to describe the evolution of the whole joint distribution of assets returns.\\newline\n\\indent As previously discussed, one of the main feature of the proposed model, concerns its ability to easily incorporate exogenous information in the state dependent correlation dynamics. Consequently, to control for the general economic conditions, we use observations of the following macroeconomic regressors as suggested by \\cite{adrian_brunnermeier.2011}, \\cite{chao_etal.2012} and \\cite{bernardi_etal.2015}:\n\\begin{enumerate}\n\\item[(I)] the weekly change in the three--month Treasury Bill rate (3MTB);\n\\item[(II)] the weekly change in the slope of the yield curve (TERMSPR), measured by the difference of the 10-year Treasury rate and the 3--month Treasury Bill rate;\n\\item[(III)] the weekly change in the credit spread (CREDSPR) between 10--year BAA rated bonds and the 10-year Treasury rate.\n\\end{enumerate}\nHistorical data for all the considered exogenous regressors are freely available on the Federal Reserve Bank of St. Louis (FRED) web site.\n\\subsection{Parameters estimation and goodness--of--fit results}\n\\label{sec:parameters_est}\n\\noindent The empirical experiment is conducted as follows. First, we estimate the marginals and the MS copula DCC models over the in--sample period, and then we make a one--step ahead rolling forecast of the multivariate joint density of the returns for the entire out of sample period. During the rolling forecast exercise, the rational investor takes an investment decision in terms of wealth allocation, at each point in time $t+s$ conditional to the information set available at time $t+s-1$, for $s=1,2,\\dots,448$. For both the marginal and copula models we re--estimate the model parameters each 24 observations, corresponding to six months using a fixed moving window.\\newline\n\\indent Table \\ref{tab:Marginal_estimates} reports the estimated coefficients for each GAS--AST marginal model. We find the scores coefficients associated with the scale $\\left(\\sigma\\right)$ and shape $\\left(\\nu\\right)$ parameters strongly significant indicating that our marginal model is able to capture the changes in the volatility patterns as well as in the tails of the conditional marginal distributions. On the contrary, the scores coefficients associated with the location $\\left(\\mu\\right)$ and the skewness $\\left(\\gamma\\right)$ parameters of the conditional AST distribution are not statistically different from zero. Moreover, we find that the shape and scale parameters of the AST distribution display high persistence, while the location and skewness parameters do not display such behaviour. The coefficients related to the scaled scores have quite similar magnitude across all the assets, indicating that the GAS updating mechanism adapts quite well to the different characteristics of the returns.\\newline\n\\indent Before moving to the copula specification, we asses the goodness of fit of the marginal models. In particular, we want to test if the PITs of the estimated marginal densities are independently and identically distributed uniformly on the unit interval. To this end, we perform the same test employed by \\cite{vlaar_palm.1993}, \\cite{jondeau_rockinger.2006b}, and \\cite{tay_etal.1998}. The test of the iid Uniform$\\left(0,1\\right)$ assumption consists of two parts.\nThe first part concerns the independent assumption, and it tests if all the conditional moments of the data, up to the fourth one, have been accounted for by the model, while the second part checks if the AST assumption is reliable by testing if the PITs are Uniform over the inverval $\\left(0,1\\right)$. In order to test if the PITs are independently and identically distributed, we define $\\hat u_{i,t}=F_\\mathcal{AST}\\left(y_{i,t},\\hat\\mathbf{\\theta}_{i,t}\\right)$ as the PIT of return $i$ at time $t$ and $\\bar{u}_i = T^{-1}\\sum_{t=1}^T \\hat u_{i,t}$ as the empirical average of the $i$--th PIT series, then we examine the serial correlation of $\\left(\\hat u_{i,t} - \\bar{u}_i\\right)^k$ for $k=1,2,\\dots,4$ by regressing each $\\left(\\hat u_{i,t} - \\bar{u}_i\\right)$ for $i=1,2,\\dots,N$ on its own lags up to the order 20.\nThe hypothesis of no serial correlation can be tested using a Lagrange multiplier test defined by the statistics $\\left(T-20\\right)R^2$ where $R^2$ is the coefficient of determination of the regression. This test is distributed according to a $\\chi^2\\left(20\\right)$ with a critic value of about 31.4 at the 5\\% confidence level. The first four columns of Table \\ref{tab:Uniform_test}, named ${\\rm DGT}-{\\rm AR}^{\\left(k\\right)}$, report the test for $k=1,2,3,4$. Our results confirm that, for almost all the considered time series and all the moments of the univariate conditional distributions, the marginal models are able to effectively account for the dynamic behaviour of the series. For what concerns the Uniform assumption of the PITs, we again employ the test suggested by \\cite{tay_etal.1998} which consists in splitting the empirical distribution of $\\hat u_{i,t}$ into $G$ bins and test whether the empirical and the theoretical distributions significantly differ on each bin. More precisely, let us define $n_{q,i}$ as the number of $\\hat u_{i,t}$ belonging to the bin $q$, it can be shown that\n\\begin{equation}\n\\xi_i=\\frac{\\sum_{q=1}^G\\left(n_i-\\mathbb{E} n_i\\right)^2}{\\mathbb{E} n_i} \\sim \\chi ^2\\left(G-1\\right),\n\\end{equation}\nwhere $G$ is the number of bins. In our case, since we have estimated the model parameters, the asymptotic distribution of $\\xi_i$ is bracketed between a $\\chi ^2\\left(G-1\\right)$ and $\\chi ^2\\left(G-m-1\\right)$ where $m$ is the number of estimated parameters.\nIn order to be consistent with the results reported by \\cite{tay_etal.1998} and \\cite{jondeau_rockinger.2006b}, we choose $G=20$ bins using a $\\chi^2\\left(19\\right)$ distribution for $\\xi_i,\\quad i=1,2,\\dots, N$, see also \\cite{vlaar_palm.1993}. The test statistic of the Uniform test are reported in the last column of Table \\ref{tab:Uniform_test} and are named $\\mathrm{DGT-H}\\left(20\\right)$.\nWe can observe that, for each series, the PITs are uniformly distributed over the interval $(0,1)$ at a confidence level lower then $1\\%$.\nThese findings definitely confirm the adequacy of our assumption on the innovation term and the conditional returns dynamic.\\newline\n\\indent We now move to study the dependence structure of the series using the proposed copula model. As previously discussed, various parameterisations of the proposed model for the conditional dependence structure of assets returns can be employed. Given the available data set, several different levels of flexibility are possible. Model specification concerns the selection of either the number of states, or the `` simple'', or `` generalised'', specifications of the joint model discussed in Section \\ref{sec:multivariate_model}, or the inclusion of the information coming from the exogenous data.\nIn particular, the choice of the number of regimes, in the HMM literature is an open question. As argued by \\cite{lindsay.1983} and \\cite{bohning.2000}, the number of states can be either estimated or tested. However given the large number of model combinations we are considering here, we decide to follow the most common practice of choosing the number of states according to the BIC, see e.g. \\cite{cappe_etal.2005}, \\cite{fruhwirth_schnatter.2006} and \\cite{zucchini_macdonald.2009}. We chose the best model according to the BIC since we want to penalise more the models with a high number of parameters. Moreover, since the marginal models are invariant to the choice of the possible dependence specifications, we report the log--likelihood and the information criteria only for the copula specifications.\nTable \\ref{tab:Model_Choice} reports the log--likelihood evaluated at its optimum, the BIC and the AIC for different combinations of number of states, type of dynamic and inclusion of exogenous covariates. According to the BIC, the best model is the '' simple\\qmcsp DCC dynamic specification with two regimes and the inclusion of the exogenous covariates. The BIC clearly select the model with the highest ability to represent different levels of the dependence structure affecting the data, and to react on the exogenous information depending on the state of the world. Table \\ref{tab:CopulaCoef} reports the estimated coefficients for the best fitting model.\nWe can observe that both regimes are characterised by strong persistence in the dynamic of the conditional dependence, however the dynamics react differently to the new endogenous and exogenous information. More precisely, we can observe that in the second regime, the coefficient $\\mathbf{A}^s$, $s=2$, which measures the reaction to new endogenous information, is much higher than the same coefficient in the other regime. This means that the second regime reacts more quickly, compared to the first one, when a change occurs in the market dependence structure. On the contrary, we observe a similar persistent behaviour of the correlation matrix within the two regimes. For what concerns the estimated state specific degree of freedom, we note that the second regime is characterised by a smaller coefficient indicating higher positive as well as negative tail dependence compared to the second one. These finding is also consistent with the empirical evidence reported by \\cite{chollete_etal.2009}, \\cite{bernardi_etal.2015} and \\cite{bernardi_catania.2015a}, in related studies.\nLooking at the estimated coefficients associated to the exogenous covariates, we can observe different behaviours of the conditional correlation matrix, depending on the regimes. In particular, we note that the returns dependence, conditional of being in the first regime, appears to react less to an absolute increase of the considered exogenous covariates. On the contrary, the dependence structure conditional of being in the second regime, displays a stronger reaction to the macroeconomic factors here considered.\nThe estimated transition probabilities of the unobserved Markov chain are $p_{11}=0.984$ and $p_{22}=0.987$, indicating quite strong persistence and an expected duration of the first and second states of about 61.5 and 75.9 weeks, respectively. Finally, we found that the dependence behaviour of the considered US market returns switched according to the 2000--2003 dot--com bubble and the recent 2008--2010 Global Financial Crisis.\n\\subsection{Portfolio allocation}\n\\noindent Now we consider the portfolio allocation problem. As stated in the Introduction, we focus on a rational investor who maximise his\/her expected utility. Specifically, once the models parameters have been estimated over the in sample period, the investor forecasts the one--step ahead joint distribution of his\/her assets return for each time $t$ of the out--of--sample period, and then he\/she allocates the available wealth among those assets accordingly to the portfolio allocation procedure discussed in Section \\ref{sec:asset_allocation}. We assume an initial wealth of 1\\textdollar, an allocation strategy that allows for either long and short positions and a risk free rate equal to zero. The choice of the level of the initial wealth is otherwise irrelevant because it does not influence the risk aversion, and hence, the allocation strategy under the defined CRRA utility function reported in equation \\eqref{eq:utility}. Moreover, in order to evaluate the influence of the level of risk aversion in the portfolio allocation strategy, we consider different levels of relative risk aversion, i.e., $\\upsilon=3,7,10,20$. In order to asses the performance of the proposed model from a portfolio optimisation viewpoint, we run an horse race with five common alternatives:\n\\begin{itemize}\n\\item[-] the `` Equally Weighted portfolio'', (EW) strategy, where the weights associated to each asset are kept fixed and equal to $1\/N$ during the validation period;\n\\item[-] the `` Minimum Variance portfolio'', (MV) strategy, where the weights associated to each asset are kept fixed and equal to those minimising the portfolio variance during the in--sample period;\n\\item[-] the `` DCC\\qmcsp strategy, where the one--step ahead forecast of the conditional mean and the conditional variance--covariance matrix of the joint distribution of assets return are performed using the DCC$\\left(1,1\\right)$ model of \\cite{engle.2002}, whose parameters are estimated by Quasi--ML, see, e.g., \\cite{francq_zakoian.2011}. The weights of each asset are then chosen by maximising equation \\eqref{eq:maximisation} truncated at the second order;\n\\item[-] the `` Naive Mean--Variance'', (NMV) and `` Naive Higher--Moments'', (NHM) strategies, where the optimal weights are chosen again by maximising equation \\eqref{eq:maximisation}, as in the case of the `` DCC\\qmcsp strategy, and the first four moments of the joint distribution of assets return are estimated using a rolling window approach.\n\\end{itemize}\nHereafter, we label the optimal investment strategy according to the model detailed in Section \\ref{sec:model}, Flexible Dynamic Dependence Model (FDDM).\\newline\n\\indent One easy way to compare different portfolios performance in terms of realised utility is to consider the so called `` management fee\\qmcsp approach. The management fee approach relies on the definition of the amount of money, $\\vartheta$, a rational investor is willing to pay, or gain, in order to switch from the portfolio that he\/she is currently holding on a given period, to a different one, and it coincides with the solution of the the following equation\n\\begin{equation}\nS^{-1}\\sum_{t=F+1}^{F+S}\\mathcal{U}\\left(1+\\mbox{\\boldmath $\\lambda$}^{\\mathcal{A}\\prime}_{t\\vert t+1}\\mathbf{y}_{t+1}\\right) = S^{-1}\\sum_{t=F+1}^{F+S}\\mathcal{U}\\left(1+\\mbox{\\boldmath $\\lambda$}^{\\mathcal{B}\\prime}_{t\\vert t+1}\\mathbf{y}_{t+1}-\\vartheta\\right),\n\\label{eq:managementFee}\n\\end{equation}\nwhere $F$ and $S$ are the length of the in--sample and out--of--sample periods, respectively, and, as before, $\\mbox{\\boldmath $\\lambda$}_{t\\vert t+1}$ denotes the vector of amounts of wealth $\\lambda_{j,t\\vert t+1}$, the investor is allocating to each asset $j=1,2,\\dots,N$ during the period $\\left(t,t+1\\right]$. From equation \\eqref{eq:managementFee}, it is easy to see that if $\\vartheta>0$, the investor is willing to pay a positive amount of money in order to switch from portfolio $\\mathcal{A}$ to portfolio $\\mathcal{B}$. On the contrary, if $\\vartheta<0$, the investor is going to ask for a higher return from portfolio $\\mathcal{B}$ in order to compensate the loss of utility he\/she would experience for switching from $\\mathcal{A}$ to $\\mathcal{B}$. Finally, if $\\vartheta=0$ the two portfolios give exactly the same utility to the investor, leaving he\/she indifferent between the two options. Table \\ref{tab:MF_mSR} reports the evaluated management fees for all the possible couples made by he FDDM strategy and a competitor, under different values of relative risk aversion coefficient $\\upsilon=\\left(3,5,10,20\\right)$. More precisely, looking at equation \\eqref{eq:managementFee}, we fix $\\mathcal{B}$ equal to the FDDM strategy, and we let $\\mathcal{A}$ to vary among the competing alternatives, i.e., the strategy a $\\mathcal{A}$ is selected among the alternatives $\\{\\mathrm{EW},\\mathrm{MV},\\mathrm{DCC},\\mathrm{NMV},\\mathrm{NHM}\\}$.\nAs a consequence, if the reported management fee is positive, then the optimal portfolio based on the proposed model is better compared to the specific competitor, otherwise the alternative strategy is preferred. Another useful indicator to compare portfolio strategies is the Sharpe ratio. However, since this indicator does not result in a measure of outperformance over alternative strategies with different levels of risk, we employ the so called modified Sharpe Ratio (mSR) introduced by \\cite{graham_harvey.1997} which is defined as\n\\begin{equation}\n\\mathrm{mSR}\\left(\\mathcal{A},\\mathcal{B}\\right)=\\frac{\\sigma_{\\mathcal{A}}}{\\sigma_\\mathcal{B}}\\mu_\\mathcal{B} - \\mu_\\mathcal{A},\n\\label{eq:modSharpeRatio}\n\\end{equation}\nwhere it immediately follows that if $\\mathrm{mSR}\\left(\\mathcal{A},\\mathcal{B}\\right)>0$, then model $\\mathcal{B}$ is preferred to model $\\mathcal{A}$ under a simple mean--variance preference ordering.\\newline\n\\indent Table \\ref{tab:MF_mSR} reports the management fee and the mSR for all the considered alternatives against the FDMM strategy. We note that, almost all the management fees and the Sharpe ratios are positive and significantly different from zero, indicating economic value in favour of the FDDM strategy. Indeed, concerning the evaluated management fees, we found that the FDDM strategy is almost always preferred by a rational investors with arbitrary level of risk aversion. The only exceptions are when the FDDM strategy is compared with the DCC strategy assuming an high degree of risk aversion $\\left(\\upsilon=3\\right)$ and with the EW strategy assuming a low degree of risk aversion $\\left(\\upsilon=20\\right)$. The results associated with the modified Sharpe ratio are in line with those reported by the management fees analysis. The only difference is that, when the degree of risk aversion is low, $\\left(\\upsilon=20\\right)$, the statistical significance of the modified Sharpe ratio vanishes, indicating that there is no evidence of outperformance between the FDDM and the alternative strategies.\nTable \\ref{tab:avg_weights} reports the average portfolio weights over the whole out of sample period associated with each strategy. We note that the positions taken by the FDDM and DCC strategies are quite homogeneous in terms of which asset to buy and which asset to sell. Indeed, the N225 and the FCHI indexes are always shorted by the FDDM and DCC strategies. These finding suggest that the signal coming from the two strategies is similar, demonstrating that a truly dynamic model is important when investment positions are concerned. We also note that, the weights associated with the NMV, NHM, and MVP are very similar, independently from the assumed degree of risk aversion of the investor. Looking at the descriptive statistics of the indexes given in Table \\ref{tab:Index_data_summary_stat} we can observe that the assets with the higher (in absolute value) wealth allocation, according to the FDDM strategy, are those who, ex--post, result in lower levels of kurtosis. On the contrary, the indexes with the lower exposition, report high kurtosis during the out--of--sample period. More interesting, looking again at Table \\ref{tab:Index_data_summary_stat}, we note that, during the in--sample period, the levels of kurtosis among assets were quite homogeneous, which means that our model is able to effectively forecast the changes in the behaviour of the joint density.\nMore precisely, on the one hand, the GAS updating mechanism is suited to effectively move the parameters through the direction of higher probability mass, and, on the other hand, the Markov switching dynamic copula model helps in capturing the evolution of the dependence structure during time and across states.\n\\section{Conclusion}\n\\label{sec:conclusion}\n\\noindent In this paper we propose a new flexible copula model being able to account for a wide variety of stylised facts affecting multivariate financial time series. Specifically, we allow the dependence structure affecting the considered financial time series to depend either on the realisation of a first order Markov chain and on exogenous covariates that may represent economic factors such as, interest rates, GDP growth rates, as well as financial indexes such as realised volatility indexes. The proposed solution is highly flexible, since the dependence dynamics account for either observable and unobservable latent factors.\nAnother interesting contribution concerns the marginal model specification. We develop a novel GAS filter \\citep{creal_etal.2013, harvey.2013} based on the Asymmetric Student--t distribution of \\cite{zhu_galbraith.2010} where all the location, scale and shape parameters evolves according to a first order transition dynamic based on the scaled score of the conditional distribution. In this way we allow for time variation in the first four moments of the conditional marginal distribution.\nParameters estimation is performed by exploiting the two step Inference Function for Margins procedure detailed in \\cite{patton.2006a} for conditional copulas, where in the second step we made use of the Expectation--Maximisation algorithm of \\cite{dempster_etal.1977}. Further details about the employed estimation technique can be found in \\cite{bernardi_catania.2015a}.\\newline\n\\indent The proposed model is applied to predict the evolution of the weekly returns of five major international stock indexes, previously considered by \\cite{jondeau_rockinger.2006a}, over the period 1988--2015. In the empirical application, in sample and out of sample performances of the proposed model are deeply investigated. The in sample analysis, reveals statistical evidence for nonlinearities and highlights the relevance of the inclusion of exogenous covariates in explaining the dependence structure affecting the indexes returns. Furthermore, our analysis confirms that our model can be effectively used to describe and predict various aspects of the conditional joint distribution of equity returns, such as time varying means, variances, co--skewness and co--kurtosis. Out of sample model performances are evaluated by solving an optimisation problem where a rational investor with power utility function takes future investment positions according to the predictive density of assets returns.\nSpecifically, we compare the optimal portfolio strategy implied by the proposed model with that of several static and dynamic alternatives. The portfolio optimisation strategy is based on the Taylor expansion of a CRRA utility function of the rational investor and involves the evaluation of the portfolio moments up to the fourth order. Since higher order moments of the joint distribution are not analytically known we resort to numerical integration. Our empirical results suggest that the proposed model outperforms competitors in terms of economic value.\n\\section*{Acknowledgments}\n\\noindent This research is supported by the Italian Ministry of Research PRIN 2013--2015, ``Multivariate Statistical Methods for Risk Assessment'' (MISURA), and by the ``Carlo Giannini Research Fellowship'', the ``Centro Interuniversitario di Econometria'' (CIdE) and ``UniCredit Foundation''. We would like to express our sincere thanks to all the participants to the Bozen -- Risk School, Bozen, September 2015, for their constructive comments that greatly contributed to improving the final version of the paper. A special thank goes also to Juri Marcucci and to all the participants of the 2015 Workshop in Econometrics and Empirical Economics (WEEE) held at the SADiBA, Perugia 27--28 August 2015, organised by the CIdE and the Bank of Italy.\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nLarge number of inflationary cosmology models is based on scalar fields, which \nplay an important role in the particle physics theories. The inflation is \nproduced by an homogeneous scalar field, dubbed inflaton, which under suitable \nconditions may lead to an early-time accelerated expansion. Following the first \nproposal of Guth~\\cite{Guth} and Sato~\\cite{Sato}, in the last years many \ninflationary models based on scalar fields ( and inspired by modified gravity \ntheories, string theories, quantum effects in the hot universe, etc) have been \nproposed.\n\nTypically the magnitude of scalar field is very large at the beginning of the \ninflation and\nthen it rolls down towards a potential minimum where the inflation ends (see \nRef.~\\cite{chaotic} as an example of chaotic inflation). In other models the \nfield can fall in a potential hole where it starts to oscillate and the \nrehating processes take place~\\cite{buca1, buca2, buca3, buca4}.\nSome more complicated models are based on a phase transition between two scalar \nfields: they are the so-called hybrid or double inflation models~\\cite{ibrida1, \nibrida2}.\nFor the introduction to the dynamics of inflation see Ref.~\\cite{Linde} and \nRefs.~\\cite{revinflazione}--\\cite{rr2}.\n\nRecently, cosmological and astrophysical data~\\cite{cosmdata} seem to confirm \nthe predictions of Starobinsky inflationary model~\\cite{Staro}. Such a model is \nbased on the account of $R^2$-term as the correction in the Einstein equations. \nThis quadratic correction emerges in the Planck epoch and plays a fundamental \nrole in the high curvature limit, when the early-time acceleration takes place. \nSuch theory is conformally equivalent to a scalar-tensor theory in the Einstein \nframe, where the inflaton drives the expansion in a quasi-de Sitter space-time \nand slowly moves to the end of inflation, when the reheating \nprocesses~\\cite{r1,r2,r3} start.\nSuch inflationary model has been recently revisited in many works. Among them, \nin Ref.~\\cite{genStar} a superconformal generalization of such a model in \nsuperconformal theory has been investigated, and in \nRefs.~\\cite{SS1}--\\cite{SS2} other applications based on the spontaneous \nbreaking of conformal invariance\nand on the scale-invariant extensions of Starobinsky model have been presented.\nIn Ref.~\\cite{cinesi}, a generalization of the Starobinsky model represented by \na polynomial correction of the Einstein gravity of the type $c_1 R^2+c_2 R^n$ \nhas been studied.\n\nIn this paper, we will concentrate on inflation caused by scalar curvature \ncorrections to Einstein gravity, namely,\nwe will consider the so-called $F(R)$-gravity, whose action is in the form of \n$F(R)=R+f(R)$, being $f(R)$ a function of the Ricci scalar (for recent reviews \non modified gravity,\nsee Refs.~\\cite{Review-Nojiri-Odintsov} --\\cite{SebRev} and \nRef.\\cite{others}). This kind of corrections may occur due to quantum effects \nin the\nhot universe or maybe motivated by the ultraviolet completation of quantum \ntheory of gravity. Our aim is to investigate which kinds of viable inflation \ncan be realized in the contest of $F(R)$-gravity beyond the Starobinsky model, \nwhose dynamics is governed in the Einstein frame by a potential of the type \n$V(\\sigma)\\sim (1-\\exp[-\\sigma])^2$, being $\\sigma$ the inflaton.\n\nThe paper is organized in the following way. In Sections {\\bf 2}--{\\bf 3}, we \nwill revisit the conformal transformations which permit to pass from the Jordan \nframe to the Einstein one and we will recall the dynamics of the viable \ninflation. In Section {\\bf 4}, we will study inflation for the general class of \nscalar potentials of the type $V(\\sigma)\\sim\\exp[n\\sigma]$, $n$ being a general \nparameter,\nperforming the analysis in the Einstein frame and therefore reconstructing the \n$F(R)$-gravity theories which correspond to the given potentials. Viable \ninflation must be consistent with the last Planck data (spectral index, \ntensor-to-scalar ratio...) and must correspond to Einstein theory corrections \nwhich emerge only at mass scales larger than the Planck one, namely at high \ncurvatures. We will see the conditions on $n$ for which slow-roll conditions \nare satisfied, and we will reconstruct the form of the $F(R)$-models during \nearly-time acceleration and their form at small curvaures. Some specific \nexamples are presented.\nIn Section {\\bf 5}, following the recent success of\nhigher derivative gravity, we will revisit and study in detail the specific \nclass of models $F(R)=R+(R+R_0)^n$. The analysis in the\nEinstein frame reveals that $n$ must be very close to two in order to realize a \nviable inflation for large and negative values of the scalar field,\nbut other possibilities are allowed\nby bounding the field in a different way.\nIn particular, when $n>2$, the de Sitter solution emerges, but in order to \nstudy the exit from inflation is necessary to analyze the theory in the Jordan \nframe, where perturbations make possible an early time acceleration with a \nsufficient amount of inflation.\nSome summary and outlook are given in Section {\\bf 6}.\nTechnical details and further considerations are presented in Appendixes.\n\n\n\\paragraph*{} We shall use units in which $c=\\hbar=k_{\\mathrm{B}}=1$,\n$c\\,,\\hbar\\,,k_{\\mathrm{B}}$ being respectively the speed of light, the\nPlanck and Boltzmann constants. Moreover we shall denote by $G_N$\nthe gravitational constant and by $M_{\\mathrm{Pl}} =G_{N}^{-1\/2} =1.2\n\\times 10^{19}$GeV the Planck mass.\nFinally, we shall set $\\kappa^2\\equiv 8 \\pi G_{N}$.\n\n\n\n\\section{Conformal transformations}\n\nIn scalar-tensor theories of gravity, a scalar field coupled to the metric \nappears in the action.\nThe first scalar-tensor theory was proposed by Brans \\& Dicke in \n1961~\\cite{BransDicke} in the attempt to incorporate the Mach's principle into \nthe theory of gravity, but today the interest to such theories is related with \nthe possibility to reproduce the primordial acceleration of the inflationary \nuniverse.\n\nIn principle, a modified gravity theory can be rewritten in scalar-tensor or \nEinstein frame\nform.\nLet us start by considering the general action of $F(R)$-modified gravity\n\\begin{equation}\nI = \\int_\\mathcal{M} d^4 x \\sqrt{-g} \\left[ \\frac{F(R)}{2\\kappa^2}\\right]\\,,\n\\label{action}\n\\end{equation}\nwhere $F(R)$ is a function of the Ricci scalar $R$, $g$ is the determinant of \nthe metric tensor $g_{\\mu\\nu}$ and\n$\\mathcal{M}$ is the space-time manifold.\nNow we introduce the field $A$ into\n~(\\ref{action}),\\\\\n\\phantom{line}\n\\begin{equation}\nI_{JF}=\\frac{1}{2\\kappa^{2}}\\int_{\\mathcal{M}}\\sqrt{-g}\\left[\nF_A(A) \\, (R-A)+F(A)\\right] d^{4}x\\label{JordanFrame}\\,.\n\\end{equation}\n\\phantom{line}\\\\\nHere, `$JF$' means `Jordan frame', and $F_A(A)$ denotes the derivative of \n$F(A)$ with respect to $A$. By making the variation\nwith respect to $A$, we immediatly obtain $A=R$, such that (\\ref{JordanFrame}) \nis equivalent to (\\ref{action}) .\nWe define the scalar field $\\sigma$ (which in fact encodes the new degree of \nfreedom in the theory, namely the scalaron or inflaton) as\n\\begin{equation}\n\\sigma := -\\sqrt{\\frac{3}{2\\kappa^2}}\\ln [F_A(A)]\\,.\\label{sigma}\n\\end{equation}\nBy considering the conformal transformation of the metric,\n\\begin{equation}\n\\tilde g_{\\mu\\nu}=\\mathrm{e}^{-\\sigma}g_{\\mu\\nu}\\label{conforme}\\,,\n\\end{equation}\nwe finally get the `Einstein frame' action\n\\phantom{line}\n\\begin{eqnarray}\nI_{EF} &=& \\int_{\\mathcal{M}} d^4 x \\sqrt{-\\tilde{g}} \\left\\{ \n\\frac{\\tilde{R}}{2\\kappa^2} -\n\\frac{1}{2}\\left(\\frac{F_{AA}(A)}{F_A(A)}\\right)^2\n\\tilde{g}^{\\mu\\nu}\\partial_\\mu A \\partial_\\nu A - \n\\frac{1}{2\\kappa^2}\\left[\\frac{A}{F_A(A)}\n- \\frac{F(A)}{F_A(A)^2}\\right]\\right\\} \\nonumber\\\\ \\nonumber\\\\\n&=&\\int_{\\mathcal{M}} d^4 x \\sqrt{-\\tilde{g}} \\left( \n\\frac{\\tilde{R}}{2\\kappa^2} -\n\\frac{1}{2}\\tilde{g}^{\\mu\\nu}\n\\partial_\\mu \\sigma \\partial_\\nu \\sigma - \nV(\\sigma)\\right)\\,,\\label{EinsteinFrame}\n\\end{eqnarray}\n\\phantom{line}\\\\\nwhere $\\tilde{R}$ denotes the Ricci scalar evaluated in the conformal metric \n$\\tilde{g}_{\\mu\\nu}$ and $\\tilde{g}$ is the determinant of the conformal \nmetric, namely $\\tilde{g}=\\mathrm{e}^{-4\\sigma}g$.\nFurthermore, one has\\\\\n\\phantom{line}\n\\begin{equation}\nV(\\sigma)\\equiv\\frac{A}{F'(A)} - \n\\frac{F(A)}{F'(A)^2}=\\frac{1}{2\\kappa^2}\\left\\{\\mathrm{e}^{\\left(\\sqrt{2\\kappa^2\/3}\\right)\\sigma}\nR\\left(\\mathrm{e}^{-\\left(\\sqrt{2\\kappa^2\/3}\\right)\\sigma}\\right)\n-\\mathrm{e}^{2\\left(\\sqrt{2\\kappa^2\/3}\\right) \n\\sigma}F\\left[R\\left(\\mathrm{e}^{-\\left(\\sqrt{2\\kappa^2\/3}\\right)\\sigma}\\right)\\right]\n\\right\\}\\,,\\label{V}\n\\end{equation}\n\\phantom{line}\\\\\nbeing $R(\\mathrm{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma})$ the solution of\nEq.~(\\ref{sigma}) with $A=R$, namely $R$ a function of \n$\\mathrm{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma}$.\nIn what follows, we will omit the tilde to denote all the quantities evaluated \nin the Einstein frame.\n\nNote that string-inspired inflationary models also contain canonical \nand\/or tachyon scalar (for recent discussion see Refs.~\\cite{Samiinfl}--\\cite{Jamil:2013nca}). \n\n\n\\section{Dynamics of inflation}\n\nIn this Section, for the sake of completeness, we will recall the well-known \nfacts on inflation. The energy density and pressure of the inflaton $\\sigma$ \nare given by\n\\begin{equation}\n\\rho_\\sigma=\\frac{\\dot\\sigma^2}{2}+V(\\sigma)\\,,\\quad \np_\\sigma=\\frac{\\dot\\sigma^2}{2}-V(\\sigma)\\,,\n\\end{equation}\nwhere the dot is the derivative with respect to the cosmological time. The \nFriedmann equations in the presence of $\\sigma$ read\n\\begin{equation}\n\\frac{3 \nH^2}{\\kappa^2}=\\frac{\\dot\\sigma^2}{2}+V(\\sigma)\\,,\\quad-\\frac{1}{\\kappa^2}\\left(2\\dot \nH+3 H^2\\right)=\\frac{\\dot\\sigma^2}{2}-V(\\sigma)\\,,\n\\end{equation}\nand the energy conservation law coincides with the equation of motion for \n$\\sigma$ and reads\n\\begin{equation}\n\\ddot\\sigma+3H\\dot\\sigma=-V'(\\sigma)\\,,\n\\end{equation}\nwhere the prime denotes the derivative of the potential with respect to \n$\\sigma$.\nFrom the Friedmann equations we obtain\n\\begin{equation}\n\\dot{H}=-\\frac{\\kappa^2}{2}\\left( \\rho_\\sigma+ p_\\sigma \n\\right)=-\\frac{\\kappa^2}{2} \\dot\\sigma^2\\,.\n\\end{equation}\nOn the other hand, the acceleration can be expressed as\n\\begin{equation}\n\\frac{\\ddot{a}}{a}= H^2+\\dot{H}=H^2\\left(1-\\epsilon \\right)\\,,\n\\end{equation}\nwhere we have introduced the ``slow roll'' parameter\n\\begin{equation}\n\\epsilon=-\\frac{\\dot{H}}{H^2}\\,.\n\\end{equation}\nThis parameter may be expressed as a function of the inflaton as\n\\begin{equation}\n\\epsilon=\\frac{\\kappa^2 \\dot\\sigma^2}{2 H^2} \\,.\n\\end{equation}\nWe also have\n\\begin{equation}\n\\frac{\\ddot{a}}{a}=\\frac{\\kappa^2}{3}\\left(V-\\dot\\sigma^2 \\right)\\,.\n\\end{equation}\nThus, the condition to have an acceleration is $\\epsilon <1$ or $\\dot\\sigma^2 \n< V(\\sigma)$. There is another slow roll parameter defined by\n\\begin{equation}\n\\eta=-\\frac{\\ddot{H}}{2H\\dot{H}}=\\epsilon-\\frac{1}{2 \\epsilon H}\\dot\\epsilon \n\\,.\n\\end{equation}\nAs a function of the inflaton one has\n\\begin{equation}\n\\eta=-\\frac{\\ddot\\sigma}{H\\dot\\sigma} \\,.\n\\end{equation}\nFor the inflation to occur and persist for a convenient amount of time, a quasi \nde Sitter space is required, namely $\\dot H$ has to be very small, and, as a \nresult, also the two slow roll parameters have to be very small, and one has\n\\begin{equation}\n\\dot\\sigma^2\\ll V(\\sigma)\\,,\n\\end{equation}\nnamely the kinetic energy of the field has to be small during the inflation. As \na result, the Friedmann equations reduce to\n\\begin{equation}\n\\frac{3H^2}{\\kappa^2}\\simeq V(\\sigma)\\,,\n\\quad\n3H\\dot\\sigma\\simeq -V'(\\sigma)\\,.\\label{EOMs}\n\\end{equation}\nIt is easy to show that within this slow roll regime,\n the slow roll parameters may be expressed as function of the inflaton \npotential as\n\\begin{equation}\n\\epsilon=\\frac{1}{2\\kappa^2}\\left(\\frac{V'(\\sigma)}{V(\\sigma)}\\right)^2\\,,\\quad \n\\eta=\\frac{1}{\\kappa^2}\\left(\\frac{V''(\\sigma)}{V(\\sigma)}\\right)\\,.\\label{slow}\n\\end{equation}\nInflation ends when $\\epsilon\\,,|\\eta|\\sim 1$. A useful quantity which \ndescribes the amount of inflation is e-foldings number $N$ defined by\n\\begin{equation}\nN\\equiv\\ln \\frac{a_f}{a_i}=\\int_{t_i}^{t_e} H dt\\simeq \n\\kappa^2\\int^{\\sigma_i}_{\\sigma_e}\\frac{V(\\sigma)}{V'(\\sigma)}d\\sigma\\,,\\label{N}\n\\end{equation}\nwhere the indices $i,f$ are referred to the quantities at the beginning and the \nend of inflation, respectively. The required e-foldings number for inflation is \nat least $N\\simeq 60$.\nThe amplitude of the primordial scalar power spectrum is\n\\begin{equation}\n\\Delta_{\\mathcal R}^2=\\frac{\\kappa^4 V}{24\\pi^2\\epsilon}\\,,\\label{spectrum}\n\\end{equation}\nand for slow roll inflation the spectral index $n_s$ and the tensor-to-scalar \nratio are given by\n\\begin{equation}\nn_s=1-6\\epsilon+2\\eta\\,,\\quad r=16\\epsilon\\,.\\label{indexes}\n\\end{equation}\nThe last Planck data constrain these quantities as\n\\begin{equation}\nn_s=0.9603\\pm0.0073\\,,\\quad r<0.11\\,.\\label{data}\n\\end{equation}\n\n\n\\section{Reconstruction of $F(R)$ theory from the scalar potential and analysis \nof the inflation in the Einstein frame}\n\nIn this Section, we\nwill study some classes of scalar potential which produce inflation in the \nEinstein frame. The aim is to generalize the Starobinsky model\nby considering different behaviour of the scalar potential as \n$V(\\sigma)\\sim\\exp[n\\sqrt{2\\kappa^2\/3}\\sigma]$, where $n$ is the parameter on \nwhich it depends the dynamics of the inflation (slow roll parameters, spectral \nindices...). This analysis is motivated by the possibility to reconstruct \nsuitable $F(R)$ corrections to General Relativity in the corresponding Jordan \nframe. By `suitable' we mean corrections that vanish at mass scales smaller \nthan the Planck mass $M_{\\text{Pl}}$ and give rise to corrections only in the high \ncurvature limits, namely during the inflationary period.\nEvery scalar potential will be confronted with cosmological data.\n\nIn order to reconstruct the $F(R)$-gravity which corresponds to a given \npotential, we may start from Eq.~(\\ref{V}). By dividing such equation to \n$\\exp\\left[2\\sqrt{2\\kappa^2\/3}\\right]$, and then by taking the derivative with respect to \n$R$, we get\n\\begin{equation}\nR \nF_R(R)=-2\\kappa^2\\sqrt{\\frac{3}{2\\kappa^2}}\\frac{d}{d\\sigma}\\left(\\frac{V(\\sigma)}{e^{2\\left(\\sqrt{2\\kappa^2\/3}\\right)\\sigma}}\\right)\\,.\\label{start}\n\\end{equation}\nAs a result, giving the explicit form of the potential $V(\\sigma),$ thanks to \nthe relation (\\ref{sigma}), we obtain an equation for $F_R(R)$, and therefore \nthe $F(R)$-gravity model in the Jordan frame. In this process, one introduces \nthe integration constant, which has to be fixed by\nrequiring that Eq.~(\\ref{V}) holds true.\n\n\\subsection{$V(\\sigma)\\sim c_0+c_1\\exp[\\sigma]+c_2\\exp[2\\sigma]$: $R+ \nR^2+\\Lambda $-models\\label{4.1}}\n\nLet us start with the following inflationary potential\n\\begin{equation}\nV(\\sigma)=\\left[c_0+c_1\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma}+c_2\\text{e}^{2\\sqrt{2\\kappa^2\/3}\\sigma}\n\\right]\\,.\\label{cucu}\n\\end{equation}\nThis is the simplest example and is the minimal generalization of the \nStarobinsky model.\nEquation (\\ref{start}) gives\n\\begin{equation}\n2c_0F_R^2+c_1F_R-RF_R=0\\,.\\label{req}\n\\end{equation}\nAssuming $F_R \\neq 0$, one gets\n\\begin{equation}\nF_R(R)=-\\frac{c_1}{2c_0}+\\frac{R}{2c_0}\\,.\\label{FF}\n\\end{equation}\nThus, the corresponding Lagrangian of $F(R)$ gravity is given by\n\\begin{equation}\nF(R)=-\\frac{c_1}{2c_0}R+\\frac{R^2}{4c_0}+\\Lambda\\,.\\label{Staro2}\n\\end{equation}\nHere, $\\Lambda$ is a constant of integration, which can be determined by \nequation (\\ref{V}). The result is\n\\begin{equation}\n\\Lambda=\\frac{c_1^2}{4c_0}-c_2\\,.\\label{L}\n\\end{equation}\nAn important remark is in order. In order to have the correct Einstein-Hilbert \nterm, we must put $-c_1\/(2c_0)=1$. Thus, we have the class \nof modified quadratic models depending on two constants\n\\begin{equation}\nF(R)=R+\\frac{R^2}{4c_0}+c_0-c_2\\,.\\label{Staro3}\n\\end{equation}\nFurthermore, in the specific case $c_2=0$, one has\n \\begin{equation}\nF(R)=R+\\frac{R^2}{4c_0}+ c_0\\,.\\label{cogno}\n\\end{equation}\nThis is an interesting model, and in the Appendix B, we will study its static \nsperically symmetric solutions.\n\nThe other interesting case is the vanishing of cosmological constant, namely\n\\begin{equation}\nc_2=c_0.\\label{L1}\n\\end{equation}\nSince we are assuming $c_1\/(2c_0)=1$, it also follows \n$c_2=-c_1\/2$. As a consequence, we recover the\nStarobinsky model $F(R)=R+\\frac{R^2}{4c_0}$, and related Einstein frame \npotential\n\\begin{equation}\nV(\\sigma)=c_0 \\left(1-\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma}\n\\right)^2\\,.\\label{VStaro}\n\\end{equation}\nThe model (\\ref{Staro3}) is the extension of the Starobinsky model to the case \nwith cosmological constant and gives a viable inflation.\nIn what follows, we denote\n\\begin{equation}\nc_0=\\frac{\\gamma}{4\\kappa^2}\\,.\n\\end{equation}\nThe initial value of the inflaton is large (and negative) and it rolls down \ntoward the potential minimum at $\\sigma\\rightarrow 0^-$, \n$V(0)=-\\gamma\/(4\\kappa^2)<0$. During inflation ($\\sigma\\rightarrow-\\infty$) \nequations (\\ref{EOMs}) lead to\n\\begin{equation}\nH^2\\simeq\\frac{\\gamma}{12}\\,,\n\\quad\n3H\\dot\\sigma\\simeq \n\\left(\\frac{\\gamma}{\\sqrt{6\\kappa^2}}\\right)\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma}\\,.\\label{copy}\n\\end{equation}\nIt means, that a quasi de Sitter solution can be realized.\nWe must require\n\\begin{equation}\n\\gamma\\propto M^2\\,,\\quad M\\ll M_P\\,,\\label{gammacond}\n\\end{equation}\nwhere $M_P$ is the Planck mass. As a consequence, the corrections of Eintein\ngravity emerge at high curvature and one gets the accelerated expansion with \n$H\\propto\\sqrt{\\gamma}$.\nThe scalar field behaves as\n\\begin{equation}\n\\sigma\\simeq \n-\\sqrt{\\frac{3}{2\\kappa^2}}\\ln\\left[\\frac{1}{3}\\sqrt{\\frac{2\\gamma}{3}}(t_0-t)\\right]\\,,\n\\end{equation}\nwhere $t_0$ is bounded at the beginning of the inflation. If at this time \n$|\\sigma|$ is very large, the slow roll paramters (\\ref{slow}),\n\\begin{equation}\n\\epsilon=\\frac{4}{3}\\frac{1}{(2-\\text{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma})^2}\\simeq \n0\\,,\\quad|\\eta|=\\frac{4}{3}\\frac{1}{|2-\\text{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma}|}\\simeq \n0\\,,\\label{star}\n\\end{equation}\nare very small and the field moves slowly.\nThe inflation ends when such parameters are of the order of unit, namely at\n$\\sigma_e\\simeq-0.17\\sqrt{3\/(2\\kappa^2)}$.\nThe e-foldings number can be evaluated from (\\ref{N}) and reads \n($\\sigma_i\\gg\\sigma_e$),\n\\begin{equation}\nN\\simeq\\frac{3\\text{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma}}{4}\\Big\\vert^{\\sigma_i}_{\\sigma_e}\\simeq\\frac{1}{4}\\sqrt{\\frac{2\\gamma}{3}}t_0\\,.\n\\end{equation}\nThe inflation ends at $t_e=t_0-3\\sqrt{3\/(2\\alpha)}\\exp[0.17]$. For example, in \norder to obtain $N=60$, we must require \n$\\sigma_i\\simeq-\\sqrt{3\/(2\\kappa^2)}4.38\\simeq 1.07 M_{pl}$. Moreover, we may \nexpress the e-foldings numbers as\n\\begin{equation}\n\\epsilon\\simeq\\frac{3}{4N^2}\\,,\\quad|\\eta|\\simeq\\frac{1}{N}\\,.\n\\end{equation}\nThe amplitude of primordial power spectrum (\\ref{spectrum}) is\n\\begin{equation}\n\\Delta_{\\mathcal R}^2\\simeq\\frac{\\kappa^2\\gamma \nN^2}{72\\pi^2}\\propto\\frac{\\kappa^2 M^2 N^2}{72\\pi^2}\\ll \\frac{\\kappa^2M_P^2 \nN^2}{72\\pi^2}\\,,\n\\end{equation}\nand the indexes (\\ref{indexes}) result to be\n\\begin{equation}\nn_s\\simeq1-\\frac{2}{N}\\,,\\quad r\\simeq\\frac{12}{N^2}\\,.\n\\end{equation}\nSince we have $n_s>1-\\sqrt{0.11\/3}\\simeq 0.809$ when $r<0.11\\,,n_s<1$, we see \nthat these indices are compatible with (\\ref{data}). For example, for $N=60$, \none has $n_s=0.967$ and $r=0.003$.\nWe stress that this behaviour is the one of Starobinsky model, with a different \nminimum of the potential and a cosmological constant in the Jordan frame. The \nappeareance of a cosmological constant at large curvature needs some \nexplanation. It can be originated by some quantum effects or may be supported \nby a modified gravity term which makes it to vanish at small curvatures (see \nfor examples the so called `two step-models' in \nRefs.~\\cite{twostep1}--\\cite{twostep2}). However, if the cosmological constant \nis set equal to zero (or, if necessary, is set equal to an other value), the feature of the model in Einstein frame does not \nchange during the inflation: we will see that this double possibility, namely, \ntaking a cosmological constant in the Jordan frame, or taking an additional \nterm $\\propto\\exp[2\\sqrt{2\\kappa^2\/3}\\sigma]$ in the Einstein frame, does not \nmodify the proprieties of the scalar potentials during the early time \nacceleration.\n\n\n\\subsection{$V(\\sigma)\\sim\\gamma\\exp[-n\\sigma]\\,,n>0$: $c_0 \nR^{\\frac{n+2}{n+1}}$-models\\label{case}}\n\nAs a second example, we consider the following potential\n\\begin{equation}\nV(\\sigma)=\\frac{\\alpha}{\\kappa^2}\\left(1-\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma}\\right)+\\frac{\\gamma}{\\kappa^2}\\text{e}^{-n\\sqrt{2\\kappa^2\/3}\\sigma}\\,,\n\\end{equation}\nbeing $\\gamma\\,,n>0$ constants. This potential possesses a minimum in which the \nscalar field may fall at the end of inflation. Since for large and negative \nvalues of the scalar field the potential is not flat, we do not expect a de \nSitter universe, but if the slow roll conditions are\nsatisfied, we can obtain an acceleration with a sufficient amount of inflation.\n\nBy using our reconstruction,\nfrom (\\ref{start}) one derives\n\\begin{equation}\nF_R(R)-\\frac{1}{2}+\\frac{\\gamma}{2\\alpha}(2+n) \nF_R(R)^{1+n}=\\frac{R}{4\\alpha}\\,.\\label{pippo}\n\\end{equation}\nIn principle, this equation admitts many solutions. At the perturbative level, \nit is easy to find that, by putting\n\\begin{equation}\n\\alpha=-\\gamma(n+2)\\,,\\label{alphacond}\n\\end{equation}\nwe obtain\n\\begin{equation}\nF_R(R\\ll\\gamma)\\simeq 1+c_1 R+c_2 R^2+c_3 R^3+...\\,,\\label{funzioncina}\n\\end{equation}\nwhen $|c_1 R|\\,,|c_2 R^2|\\,, |c_3 R^3|...\\ll 1$ (see Appendix A). It means \nthat, since \n$c_1\\propto\\gamma^{-1}\\,,c_2\\propto\\gamma^{-2}\\,,c_3\\propto\\gamma^{-3}...$,\nif $\\gamma$ satisfies (\\ref{gammacond}), we recover the Einstein's gravity when \n$R\\ll\\gamma$ and the theory is an high curvature correction to General \nRelativity. Moreover, when $R\\gg\\gamma$, the asymptotic solution of \nEq.~(\\ref{pippo}) with (\\ref{alphacond}) is given by\n\\begin{equation}\nF_R(R\\gg\\gamma)\\simeq\\left(\\frac{1}{4(n+2)}\\right)^\\frac{1}{1+n}\\left(\\frac{R}{\\gamma}\\right)^{\\frac{1}{n+1}}\\,,\\quad\nF(R\\gg\\gamma)\\simeq\\gamma\\left(\\frac{n+1}{n+2}\\right)\\left(\\frac{1}{4(n+2)}\\right)^\\frac{1}{1+n}\\left(\\frac{R}{\\gamma}\\right)^{\\frac{n+2}{n+1}}\\,.\\label{expression}\n\\end{equation}\nHere, one important comment is required. In the Einstein frame, \n$R_{EF}\\sim\\gamma$ during inflation, but the corresponding curvature in the \nJordan frame is $R_{JF}\\simeq \\text{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma}R_{EF}$ (we \nmay neglect the kinetic energy of scalar field in the slow roll approximation), \nsuch that $R_{JF}\\gg\\gamma$ when $\\sigma\\rightarrow-\\infty$ and expression \n(\\ref{expression}), which is evaluated in the Jordan frame, effectively is \nvalid for inflation\n(for $n\\rightarrow 0$ we recover $F(R)\\sim R^2$ in the Jordan frame).\n\nThe potential finally reads\n\\begin{equation}\nV(\\sigma)=-\\frac{\\gamma(n+2)}{\\kappa^2}\\left(1-\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma}\\right)+\\frac{\\gamma}{\\kappa^2}\\text{e}^{-n\\sqrt{2\\kappa^2\/3}\\sigma}\\,.\n\\end{equation}\nThis potential has a minimum ($V'(\\sigma_{min})=0$) at \n$\\sigma_{min}=-\\sqrt{3\/(2\\kappa^2)}\\log[(n+2)\/n]\/(n+1)$, and one gets\n\\begin{equation}\nV(\\sigma_{min})=\\frac{\\gamma}{\\kappa^2}\\left(n^{\\frac{1}{n+1}}(n+2)^{\\frac{n}{n+1}}+\\left(\\frac{n+2}{n}\\right)^{\\frac{1}{n+1}}-(n+2)\\right)>0\\,,\\quad\\gamma\\,,n>0\\,.\n\\end{equation}\nWhen $\\sigma\\rightarrow-\\infty$ (large curvature), the potential goes to \ninfinity, and when $\\sigma\\rightarrow 0^-$, $V(\\sigma)=\\gamma\/\\kappa^2$.\nSince the slow roll parameters are given by\n\\begin{eqnarray}\n\\epsilon&=&\\frac{\\left(n-(n+2)\\text{e}^{(n+1)\\sqrt{2\\kappa^2\/3}\\sigma}\\right)^2}{3\\left(\n(n+2)\\text{e}^{n\\sqrt{2\\kappa^2\/3}\\sigma}-(n+2)\\text{e}^{(n+1)\\sqrt{2\\kappa^2\/3}\\sigma}-1\n\\right)^2}\\,,\\nonumber\\\\ \\nonumber\\\\\n|\\eta|&=&\\frac{2}{3}\\frac{n^2+(n+2)\\text{e}^{(n+1)\\sqrt{2\\kappa^2\/3}\\sigma}}{|1+(n+2)\\text{e}^{(n+1)\\sqrt{2\\kappa^2\/3}\\sigma}-\n(n+2)\\text{e}^{n\\sqrt{2\\kappa^2\/3}\\sigma}\n|}\\,,\n\\end{eqnarray}\none has $\\epsilon(\\sigma\\rightarrow-\\infty)\\simeq n^2\/3$ and\n$|\\eta(\\sigma\\rightarrow-\\infty)|\\simeq2n^2\/3$, which implies $00\\,,\\quad (n<\\sqrt{3})\n\\end{equation}\nand we have an acceleration as soon as $\\epsilon<1$. Despite to the fact that \nin this kind of models the acceleration is smaller than in the de Sitter \nuniverse, the slow roll paramters can be small enough\nto justify our slow roll approximations. A direct evaluation of the ratio of \nkinetic energy of the field and potential leads to \n$\\left(\\dot\\sigma^2\/2\\right)\/V(\\sigma)=n^2\/9$, which is much smaller than one \nwhen $n\\ll 1$. The inflation ends when the slow roll paramters are on the order \nof unit, before the minimum of the potential. Note that, by definition, since \n$V'(\\sigma_{min})=0$, one has $\\epsilon(\\sigma_{min})=0$, which corresponds to \nthe minimum of the slow roll paramter. However, before to this point, since the \nslow roll parameters behave as\n\\begin{equation}\n\\epsilon\\simeq\\frac{n^2}{3\\left(\n(n+2)\\text{e}^{n\\sqrt{2\\kappa^2\/3}\\sigma}-1\n\\right)^2}\\,,\\quad\n|\\eta|\\simeq\\frac{2}{3}\\frac{n^2}{|1-\n(n+2)\\text{e}^{n\\sqrt{2\\kappa^2\/3}\\sigma}\n|}\\,,\n\\end{equation}\nwe find that $\\epsilon\\,,|\\eta|\\simeq 1$ when \n$\\sigma\\simeq-\\sqrt{3\/(2\\kappa^2)}\\log[(6+3n)\/(3-n\\sqrt{3})]\/n$.\nFinally, from Eq.~(\\ref{N}), we get the $N$-foldings number of inflation,\n\\begin{equation}\nN\\simeq -\\frac{1}{n}\\sqrt{\\frac{3\\kappa^2}{2}}\\sigma\n\\Big\\vert^{\\sigma_i}_{\\sigma_e}\\simeq-\\frac{3}{n^2}\\ln\\left[\\frac{n^2}{3}\\sqrt{\\frac{\\gamma}{3}}t_0\\right]\\,.\n\\end{equation}\nWhen the inflation ends, the field falls in the minimum of the potential and \nstarts to oscillate.\nThe reheating process takes place. The amplitude of primordial power spectrum \n(\\ref{spectrum}) and the spectral indexes (\\ref{indexes}) can be written as\n\\begin{equation}\n\\Delta_{\\mathcal R}^2\\simeq\\frac{\\kappa^2\\gamma \n\\text{e}^{\\frac{2}{3}N\\,n^2}}{8\\pi^2n^2}\\,,\n\\quad\nn_s\\simeq1-\\frac{2n^2}{3}\\,,\\quad r\\simeq\\frac{16n^2}{3}\\,.\n\\end{equation}\nThe corrections to $n_s$ and $r$ are on the order of $\\exp\\left[-2n^2 \nN\/3\\right]\\ll 1$.\nThese indexes are compatible with (\\ref{data}) when\n\\begin{equation}\nn\\sim\\frac{1}{10}\\,,\\frac{2}{10}\\,.\n\\end{equation}\nThese are the typical values of $n$ which make the scalar potential (\\ref{V3}) \nable to reproduce a viable inflationary scenario. The result suggests that only \nthe models\nclose to $R^2$-gravity are able to produce this kind of inflation.\n\nIn general, $F_R(R)$ in Eq.~(\\ref{funzioncina}) may lead to a cosmological \nconstant proportional to $\\gamma$ in the $F(R)$-model. However, we can set it \nequal to zero, adding a suitable term in\nthe potential ($\\ref{V3}$) proportional to\n$\\exp\\left[2\\sqrt{2\\kappa^2\/3}\\sigma\\right]$, which\nchanges much slowler than\n$\\exp\\left[-n\\sqrt{2\\kappa^2\/3}\\sigma\\right]$ when $00$, as usually, constants. It follows\n\\footnote{Here, we exclude the solution with the minus sign in front of the \nsquare root which leads to an imaginary value of $\\sqrt{F_R}$ when $R$ is \nreal.}\n\\begin{equation}\nF_R(R)=\\frac{9\\gamma^2+8R\\alpha+3\\sqrt{16 \nR\\alpha\\gamma^2+9\\gamma^4}}{32\\alpha^2}\\,.\n\\end{equation}\nSince we must require $\\alpha\\,,\\gamma\\gg 1$ and at small curvature we want to \nrecover the Einstein gravity ($F_R=1$), we set\n$\\alpha=3\\gamma\/4$, $\\gamma>0$, and we get\n\\begin{equation}\nF_R(R)=\\frac{1}{2}+\\frac{1}{3\\gamma}R+\\frac{\\sqrt{3}}{6}\\sqrt{4R\/\\gamma+3}\\,.\\label{FFFF}\n\\end{equation}\nThus, from Eq.~(\\ref{V}) we obtain\n\\begin{equation}\nF(R)=\\frac{R}{2}+\\frac{R^2}{6\\gamma}+\\frac{\\sqrt{3}}{36}\\left(4R\/\\gamma+3\\right)^{3\/2}+\\frac{\\gamma}{4}\\,.\\label{zippozap}\n\\end{equation}\nHere, we stress that the conformal transformation gives for the Ricci scalar\n\\begin{equation}\nR=3\\text{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma}\\left(1+\n\\text{e}^{\\sqrt{2\\kappa^2\/3\/2}\\sigma\/2}\\right)\\,,\\,\n3\\text{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma}\\left(1-\n\\text{e}^{\\sqrt{2\\kappa^2\/3\/2}\\sigma\/2}\\right)\\,,\n\\end{equation}\nbut only the second one leads to our potential (namely, is the one wich emerges \nfrom our reconstruction). For $R\\ll\\gamma$, the model reads $F(R\\ll\\gamma)\\simeq R+\\gamma\/2$. If we want to recover the General Relativity action $F(R\\ll\\gamma)\\simeq R$, we must set the cosmological constant, namely the last term in (\\ref{zippozap}), equal to $-\\gamma\/4$: in this case\nthe scalar potential is\n\\begin{equation}\nV(\\sigma)=\\frac{3\\gamma}{4\\kappa^2}-\\frac{\\gamma}{\\kappa^2}\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma\/2}\n-\\gamma\\frac{\\text{e}^{2\\sqrt{2\\kappa^2\/3}\\sigma}}{4\\kappa^2}\\,.\\label{V2bis}\n\\end{equation}\nLet us analyze the possibility to reproduce inflation from the potential \n(\\ref{V2}), which finally reads\n\\begin{equation}\nV(\\sigma)=\\left(\\frac{3}{4\\kappa^2}-\\frac{\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma\/2}}\n{\\kappa^2}\\right)\\gamma\\,.\n\\end{equation}\nThe initial value of the inflaton is large (and negative) and it rolls down \ntoward the potential minimum at $\\sigma\\rightarrow 0^-$, \n$V(0)=-\\gamma\/(4\\kappa^2)<0$. When $\\sigma\\rightarrow-\\infty$ the EOMs in the \nslow roll limit read\n\\begin{equation}\nH^2\\simeq\\frac{\\gamma}{4}\\,,\n\\quad\n3H\\dot\\sigma\\simeq \n\\left(\\frac{\\gamma}{\\sqrt{6\\kappa^2}}\\right)\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma\/2}\\,.\n\\end{equation}\nIt means, that $\\gamma$ must satisfy condition (\\ref{gammacond}), such that the \ninflation takes place at the Plank epoch.\nThe de Sitter expansion can be realized and the field behaves as\n\\begin{equation}\n\\sigma\\simeq \n-\\sqrt{\\frac{6}{\\kappa^2}}\\ln\\left[\\frac{\\sqrt{\\gamma}}{9}(t_0-t)\\right]\\,,\n\\end{equation}\nwhere $t_0$ is bounded at the beginning of the inflation. If at this time the \nmagnitude of $\\sigma$ is very large, the slow roll parameters (\\ref{slow})\n\\begin{equation}\n\\epsilon=\\frac{4}{3}\\frac{1}{\\left(4-3\\text{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma\/2}\\right)^2}\\simeq\n0\\,,\\quad|\\eta|=\\frac{2}{3}\\frac{1}{|4-3\\text{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma\/2}|}\\simeq \n0\\,,\n\\end{equation}\nare small and the field moves slowly.\nThe inflation ends at\n$\\sigma_e\\simeq-0.12\\sqrt{3\/(2\\kappa^2)}$, when the slow roll parameters are of \nthe order of unit.\nThe e-foldings number can be evaluated from (\\ref{N}) and read,\n\\begin{equation}\nN\\simeq\\frac{9\\text{e}^{-\\sqrt{2\\kappa^2\/3}\\sigma\/2}}{2}\\Big\\vert^{\\sigma_i}_{\\sigma_e}\\simeq\\frac{\\sqrt{\\gamma}}{2}t_0\\,.\n\\end{equation}\nAs a consequence, the slow roll parameters can be written as\n\\begin{equation}\n\\epsilon\\simeq\\frac{3}{N^2}\\,,\\quad|\\eta|\\simeq\\frac{1}{N}\\,.\n\\end{equation}\nThe amplitude of primordial power spectrum (\\ref{spectrum}) and the spectral \nindexes (\\ref{indexes}) are given by\n\\begin{equation}\n\\Delta_{\\mathcal R}^2\\simeq\\frac{\\kappa^2\\gamma N^2}{96\\pi^2}\\,,\n\\quad\nn_s\\simeq1-\\frac{1}{N}\\,,\\quad r\\simeq\\frac{48}{N^2}\\,.\n\\end{equation}\nSince from these formulas we have $n_s>1-\\sqrt{0.11\/48}\\simeq 0.9521$ when \n$r<0.11\\,,n_s<1$, we see that these expressions are compatible with \n(\\ref{data}). For example, for $N=60$, one has $n_s=0.967$ and $r=0.013$.\n\nWe finish this Subsection with some considerations on the potential \n($\\ref{V2bis}$), which corresponds to the model with cosmological constant equal to $-\\gamma\/4$. \nSince the term $\\exp\\left[2\\sqrt{2\\kappa^2\/3}\\sigma\\right]$ changes slower than\n$\\exp\\left[\\sqrt{2\\kappa^2\/3}\\sigma\/2\\right]$, the dynamics of inflation is the \nsame of above. In this case, when the inflaton exits from the slow roll region, \nit falls in the minimum of the potential located at $V(0)=-\\gamma\/(8\\kappa^2)$.\n\n\n\\subsection{$V(\\sigma)\\sim\\gamma(2-n)\/2-\\gamma\\exp[n\\sigma]\\,,00\\,,\\label{alphaalpha}\n\\end{equation}\nif $\\gamma$ satisfies (\\ref{gammacond}),\nat small curvature one gets\n\\begin{equation}\nF_R(R\\ll\\gamma)\\simeq 1+c_1 R+c_2 R^2+c_3 R^3+...\\,,\\label{espansione}\n\\end{equation}\nwith \n$c_1\\propto\\gamma^{-1}\\,,c_2\\propto\\gamma^{-2}\\,,c_3\\propto\\gamma^{-3}...$, \nsuch that our theory is the high curvature correction to General Relativity \n(see Appendix A).\nFor example, in the previous Subsection we have seen an exact solution for the \ncase $n=1\/2$. Since $1\\ll \\gamma$, when $R\/\\gamma\\ll 1$ we can expand $F_R(R)$ \nin (\\ref{FFFF}) as\n\\begin{equation}\nF_R(R\\ll\\gamma)\\simeq \n1+\\frac{2}{3\\gamma}R-\\frac{R^2}{9\\gamma^2}+...\\,,\\label{primeF}\n\\end{equation}\nwhich returns to be (\\ref{espansione}) by using the coefficients in Appendix A.\nOn the other side, when $R\\gg\\gamma$, the asymptotic solution of \nEq.~(\\ref{pippo2}) with (\\ref{alphaalpha}) is given by\n\\begin{eqnarray}\nF_R(R\\gg\\gamma)&\\simeq&\\left(\\frac{R}{2\\gamma(2-n)}\\right)+\\left(\\frac{R}{2\\gamma(2-n)}\\right)^{1-n}\\,,\\nonumber\\\\\\nonumber\\\\\nF(R\\gg\\gamma)&\\simeq&\\frac{1}{2}\\left(\\frac{R^2}{2\\gamma(2-n)}\\right)+\\frac{1}{2-n}\\left(\\frac{1}{2\\gamma(2-n)}\\right)^{1-n}R^{2-n}\\,.\n\\label{megastana}\n\\end{eqnarray}\nAlso in this case, by taking the exact solution of the prevoius Section in the \nhigh curvature limit,\n\\begin{equation}\nF_R(R\\gg\\gamma)\\simeq\\frac{R}{3\\gamma}+\\sqrt{\\frac{R}{3\\gamma}}\\,,\n\\end{equation}\nwe can verify the consistence of expression (\\ref{megastana}) for $n=1\/2$.\n\nAs an other example, let us consider the case $n=1\/3$. The reconstruction leads \nto\\\\\n\\phantom{line}\\\\\n\\begin{equation}\nF_R(R)=\\frac{1}{3}+\\frac{3}{10}\\left(\\frac{R}{\\gamma}\\right)\n+\\frac{2}{3\\times \n5^{1\/3}}\\left[9\\left(\\frac{R}{\\gamma}\\right)+5\\right]\\frac{1}{\\Delta^{1\/3}}+\\frac{1}{6\\times\n5^{2\/3}}\\Delta^{1\/3}\\,,\n\\end{equation}\n\\phantom{line}\\\\\nwhere\n\\phantom{line}\n\\begin{equation}\n\\Delta=200+243\\left(\\frac{R}{\\gamma}\\right)^2+540\\left(\\frac{R}{\\gamma}\\right)+27\\sqrt{81\\left(\\frac{R}{\\gamma}\\right)^4+40\\left(\\frac{R}{\\gamma}\\right)^3}\\,.\n\\end{equation}\n\\phantom{line}\\\\\nAt small curvature, it is easy to find\n\\begin{equation}\nF_R(R\\ll\\gamma)\\simeq1+\\frac{9}{10}\\left(\\frac{R}{\\gamma}\\right)+...\\,,\n\\end{equation}\nand in the high curvature limit this model has the following structure,\n\\begin{equation}\nF_R(R\\gg\\gamma)\\simeq \n\\frac{3}{10}\\left(\\frac{R}{\\gamma}\\right)+\\left(\\frac{3}{10}\\right)^{\\frac{2}{3}}\\left(\\frac{R}{\\gamma}\\right)^{\\frac{2}{3}}\\,,\n\\end{equation}\nwhich corresponds to (\\ref{megastana}) with $n=1\/3$.\nFinally, the model (\\ref{FF}) is the limiting case of $n\\rightarrow 1$.\n\nLet us analyze this class of potentials. When the magnitude of the inflaton is \nlarge, the scalar potential (\\ref{V3}) with\n$\\alpha=\\gamma(2-n)\/2$,\n\\begin{equation}\nV(\\sigma)=\\frac{\\gamma(2-n)}{2\\kappa^2}-\\frac{\\gamma}{\\kappa^2}\\text{e}^{n\\sqrt{2\\kappa^2\/3}\\sigma}\\,,\\label{n<2}\n\\end{equation}\nbehaves as $V(\\sigma)\\simeq\\gamma(2-n)\/(2\\kappa^2)$, and the EOMs in the slow \nroll limit read\n\\begin{equation}\nH^2\\simeq\\left(\\frac{\\gamma(2-n)}{6}\\right)\\,,\n\\quad\n3H\\dot\\sigma\\simeq \n\\left(\\frac{n\\gamma\\sqrt{2}}{\\sqrt{3\\kappa^2}}\\right)\\text{e}^{n\\sqrt{2\\kappa^2\/3}\\sigma}\\,.\n\\end{equation}\nAs a consequence, the field results to be\n\\begin{equation}\n\\sigma\\simeq \n-\\sqrt{\\frac{3}{2\\kappa^2}}\\frac{1}{n}\\ln\\left[\\frac{2\\sqrt{2}}{3\\sqrt{3}}\\frac{n^2\\sqrt{\\gamma}}{\\sqrt{2-n}}(t_0-t)\\right]\\,,\n\\end{equation}\nwhere $t_0$ is bounded at the beginning of the inflation. When \n$\\sigma\\rightarrow-\\infty$, the slow roll parameters became\n\\begin{equation}\n\\epsilon=\\frac{4n^2}{3}\\frac{1}{\\left(2+(n-2)\\text{e}^{-n\\sqrt{2\\kappa^2\/3}\\sigma}\\right)^2}\\simeq\n0\\,,\n\\quad\n|\\eta|=\\frac{4n^2}{3}\\frac{1}{|2+(n-2)\\text{e}^{-n\\sqrt{2\\kappa^2\/3}\\sigma}|}\\simeq \n0\\,,\n\\end{equation}\nand are very small.\nThe inflation ends at\n$\\sigma_e\\simeq-(1\/n)\\sqrt{3\/(2\\kappa^2)}\\log\\left[(\\gamma\/\\alpha)(1-n\/\\sqrt{3})\\right]$,\n\nsuch that the slow roll parameters are of the order of unit and the field \nreaches the minimum of the potential at $V(0)= -\\gamma\/(4\\kappa^2)$.\nThe e-foldings number (\\ref{N}) is given by\n\\begin{equation}\nN\\simeq\\frac{3(2-n)\\text{e}^{-n\\sqrt{2\\kappa^2\/3}\\sigma}}{4n^2}\\Big\n\\vert^{\\sigma_i}_{\\sigma_e}\\simeq\n\\sqrt{\\frac{2(2-n)\\gamma}{6}}t_0\\,,\n\\end{equation}\nand\n\\begin{equation}\n\\epsilon\\simeq\\frac{3}{4n^2N^2}\\,,\\quad|\\eta|\\simeq\\frac{1}{N}\\,.\n\\end{equation}\nAs a consequence, the amplitude of primordial power spectrum and the spectral \nindexes read\n\\begin{equation}\n\\Delta_{\\mathcal R}^2\\simeq\\frac{\\kappa^2\\gamma(2-n)n^2 N^2}{36\\pi^2}\\,,\n\\quad\nn_s\\simeq1-\\frac{1}{N}\\,,\\quad r\\simeq\\frac{12}{n^2 N^2}\\,.\n\\end{equation}\nSince $n_s>1-(\\sqrt{0.11\/12})n\\simeq 1-(0.0957)n$ when $r<0.11\\,,n_s<1$, one \nhas that\n\\begin{equation}\n0.33861$. Since in this case Eq.~(\\ref{alphaalpha}) could lead to $\\alpha<0$ \n(when $n>2$) making the potential unable to reproduce inflation, we may \ngeneralize the potential to the following form\n\\begin{equation}\nV(\\sigma)=\\frac{\\alpha}{\\kappa^2}\\left(1-\n\\text{e}^{n\\sqrt{2\\kappa^2\/3}\\sigma}\\right)\n-\\frac{\\gamma}{\\kappa^2}\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma}\\,,\\label{V4}\n\\end{equation}\nwhere $\\alpha\\,,\\gamma>0$ and $n>1$.\nNow, in order to recover the Einstein gravity at small curvature, we have to \nput\n\\begin{equation}\n\\alpha=\\frac{\\gamma}{2(n-1)}>0\\,.\n\\end{equation}\nIn the slow roll limit ($\\sigma\\rightarrow-\\infty$) the EOMs read\n\\begin{equation}\nH^2\\simeq\\frac{\\gamma}{6(n-1)}\\,,\n\\quad\n3H\\dot\\sigma\\simeq \\left(\\frac{\\sqrt{2}\\gamma}{\\sqrt{3\\kappa^2}}\\right)\n\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma}\\,,\n\\end{equation}\nand the analysis of inflation results to be the same of $R^2$-models (with or \nwithout the cosmological constant term).\\\\\n\\\\\nHere, we can give some comments about the results of our investigation. Scalar \ninflationary potentials wich satisfy viable conditions for realistic primordial \nacceleration in the contest of large scalar curvature corrections to General \nRelativity can be classified in two classes. In the first one, the scalar \npotential behaves as \n$V(\\sigma)\\sim\\exp\\left[n\\sqrt{2\\kappa^2\/3}\\sigma\\right]\\,,n>0$, and produces \nacceleration with a sufficient amount of inflation only for $n$ very close to \nzero, namely the $F(R)$-model must be close to the Starobinsky one. In the \nsecond class, $V(\\sigma)\\sim\\alpha \n-\\gamma\\exp\\left[-n\\sqrt{2\\kappa^2\/3}\\sigma\\right]$, $n>0$, and a quasi de \nSitter solution emerges during inflation. Cosmological data are satisfied if \n$n>1\/3$. This kind of scalar potentials is originated from large scalar \ncurvature corrections of the type $F(R)\\sim c_1 R^2+c_2 R^\\zeta$, with \n$1<\\zeta<5\/3$ when $1\/31\\,,\n\\label{modellino}\n\\end{equation}\nwhere $\\alpha>0$ is a (dimensional) constant parameter and $R_0\\,,\\Lambda$ are \ntwo (cosmological) constants introduced to generalize the model. This class of \nmodels, following the success of quadratic correction, have been often analyzed \nin literature in order to reproduce the dynamics of inflation and\nrecently, in Ref.~\\cite{cinesi}, the cases of polynomial corrections added to \nthe Starobinsky model have been investigated.\nIn what follows, at first we would like to study the\ninflation in the Einstein frame given by (\\ref{modellino}) with $n\\neq2$. We \nwill see that inflation for large magnitude values of the field is realized for \n$n\\lesssim2$. However, other possibilities are allowed by considering \nintermediate values of the field. In this case, the de Sitter solution may \nemerge in the models with $n>2$, but in order to study the exit from inflation \nis necessary to analyze the theory in the Jordan frame, where perturbations \nmake possible an early time acceleration with a sufficient amount of \n$N$-foldings number: that is the aim of the second part of this Section.\n\nLet us start from the potential\n in the scalar field representation of (\\ref{modellino}),\n\\begin{eqnarray}\nV(\\sigma)&=&\\frac{1}{2\\kappa^2}\\left\\{\\left(\\frac{1}{\\alpha \nn}\\right)^{\\frac{1}{n-1}}\\left(1-\\text{e}^{\\sqrt{(2\\kappa^2\/3)}\\sigma}\\right)^{\\frac{n}{n-1}}\\text{e}^{\\frac{(n-2)}{(n-1)}\\sqrt{(2\\kappa^2\/3)}\\sigma}\\left[\\frac{n-1}{n}\\right]\\right.\\nonumber\\\\\n&&\\left.+R_0\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma}\\left(\\text{e}^{\\sqrt{2\\kappa^2\/3}\\sigma}-1\\right)\n-\\Lambda\\text{e}^{2\\sqrt{2\\kappa^2\/3}\\sigma}\\right\\}\\,.\n\\end{eqnarray}\nWe immediatly see that only if $n=2$, when $\\sigma\\rightarrow-\\infty$, one gets \n$V(\\sigma)\\sim\\text{const}$ and obtains the de Sitter solution (for $n=2$, \n$\\Lambda=R_0=0$ we recover (\\ref{VStaro}),\n$V(\\sigma)=(\\exp[\\sqrt{2\\kappa^2\/3}\\sigma]-1)^2\/(8\\alpha\\kappa^2)$).\nOn the other hand, even if $10$ very small bounded at the \nbeginning of inflation.\nThe Hubble parameter reads\n\\begin{equation}\nH=\\frac{3(n-1)^2}{(2-n)^2}\\frac{1}{(t_0+t)}\\,,\\quad\\frac{\\ddot \na}{a}=\\frac{3(n-1)^2((n+1)^2-2)}{(n-2)^4}\\frac{1}{(t_0+t)^2}>0\\,,\n\\end{equation}\nand we have an acceleration with decreasing Hubble parameter and curvature\n\\begin{equation}\nR=\\frac{3(n-1)^2\\left(3(n-1)^2-(2-n)\\right)}{(2-n)^4(t_0+t)}\\,.\n\\end{equation}\nThe $N$-foldings number of inflation is given by\n\\begin{equation}\nN\\simeq -\\frac{(n-1)}{(2-n)}\\sqrt{\\frac{3\\kappa^2}{2}}\\sigma\n\\Big\\vert^{\\sigma_i}_{\\sigma_e}\\simeq\n-3\\left(\\frac{n-1}{2-n}\\right)^2\\ln\\left[\n\\sqrt{\\frac{2}{3n}}\\frac{(2-n)^2}{6(n-1)^{3\/2}}\\left(\\frac{1}{n\\alpha}\\right)^{\\frac{1}{2(n-1)}}\nt_0\\right]\\,.\n\\end{equation}\nWhen the inflation ends, the field reaches the minimum of the potential. \nWe note that the cosmological constants introduced in the model do not play \nany role at inflationary stage, and change the behaviour of the potential only \nat the end of the inflation.\nIf \n$\\Lambda=R_0=0$, it is easy to see that $V(0)=0$, otherwise the potential \npossesses a minimum before $\\sigma=0$ (in the same way of \\S~\\ref{case}), where \nthe field falls and starts to oscillate.\n\nThe amplitude of primordial power spectrum (\\ref{spectrum}) and the spectral \nindexes (\\ref{indexes}) are\n\\begin{equation}\n\\Delta_{\\mathcal \nR}^2\\simeq\\frac{(n-1)^3\\kappa^2\\text{e}^{\\frac{2}{3}N\\frac{(n-2)^2}{(n-1)^2}}}{16\\pi^2(2-n)^2n}\\left(\\frac{1}{\\alpha\\,n}\\right)^{\\frac{1}{n-1}}\\,,\n\\quad\nn_s\\simeq1-\\frac{8(2-n)}{3(n-1)}\\,,\\quad r\\simeq\\frac{16(2-n)^2}{3(n-1)^2}\\,,\n\\end{equation}\nwhere we have taken into account that $n$ is close to two. In order to make \nthese indexes compatible with (\\ref{data}), we must require\n$ n\\sim 1.8\\,, 1.9$,\nsuch that, as a matter of fact, the quadratic corrections with $n=2$ appear to be \nthe only possibilities of the type (\\ref{modellino}) able to reproduce a viable \ninflation in the Einstein frame. However, it is interesting to extend our \ninvestigation to the Jordan frame.\n\n\\subsection{Inflation in the Jordan frame}\n\nIn the first part of this Section, we have seen how model (\\ref{modellino}) can \nproduce inflation in the Einstein frame representation. We have an early time \nacceleration followed by the end of inflation and the slow roll conditions are \nsatisfied if $n$ is very close to two.\nFor the sake of simplicity, in this Subsection, we \nwill put $\\Lambda=R_0=0$ in (\\ref{modellino}).\nLet us return to the Jordan frame.\n\nThe Friedmann equation for a generic $F(R)$-model read (in vacuum)\n\\begin{equation}\n\\left(R F_R(R)- F(R)\\right)-6H^2 F_R(R)-6H\\dot F_R(R)=0\\,,\\label{eq}\n\\end{equation}\nsuch that for (\\ref{modellino}) one derives\n\\begin{equation}\n\\alpha(R)^{n-1}\\left[R(1-n)\\right]-6H^2\\left[1+\\alpha \nn(R)^{n-1}\\right]=6H\\alpha n\\,(n-1)(R)^{n-2}\\dot R\\,.\\label{First}\n\\end{equation}\nIn the high curvature limit (it means, for large and negative values of the \nscalar field $\\sigma$ in the Einstein frame) we may consider $\\alpha \nn(R)^{n-1}\\gg 1$,\n\\begin{equation}\nH^2\\simeq\\frac{(n-1)}{6n}\\left[R-\\frac{6n\\,H\\,\\dot R}{R}\\right]\\,.\n\\end{equation}\nDuring inflation, the Hubble parameter $H$ evolves slowly and we can consider \nthe anologous of the slow roll parameterers which have the same\nformal dependence of the Einstein frame, namely we require $|\\dot H|\/H^2\\ll \n1$, $|\\ddot H\/(H\\dot H)|\\ll 1$. Thus, we get\n\\begin{equation}\n\\frac{\\dot \nH}{H^2}\\simeq-\\left[\\frac{4-2n}{n+4(n-1)^2}\\right]\\,,\\quad\\frac{\\ddot a}{a}=H^2\n+\\dot H=1-\\frac{4-2n}{n+4(n-1)^2}\\,.\n\\end{equation}\nIt follows that the expansion is accelerated only if\n\\begin{equation}\nn>\\frac{5}{4}=1.25\\,.\n\\end{equation}\nThis value is close to the one given by the left side of (\\ref{nn}). Moreover, \nfor $n>2$, we see that $\\dot H>0$ and the Hubble parameter, and therefore the \ncurvature, grows up with the time and the physics of Standard Model is not \nreached. This result is in agreement with the behaviour of the model discussed \nin the Einstein frame: when $n>2$ and the acceleration emerges in high \ncurvature limit, the field cannot reach a minimum of the potential for small \nvalues and we do not exit from inflation. We could arrive to this conclusion \nalso by looking for Eq.~(\\ref{bimbo}): when $n>2$, $t_0$ has to be very large \nat the beginning of inflation and the scalar field grows up with the time. \nHowever, it does not mean that models with $n>2$ do not produce inflation. In \nscalar field representation we start at very large and negative values of the \nfield (it means, that in Jordan framework we can ignore the $R$-term in the \naction). It may be interesting to see what happen for intermediate values of \nthe field. We analyze the problem in the Jordan frame. From Eq.~(\\ref{eq}) we \nhave the de Sitter condition at $R_0=12H_0^2$, $H_0$ constant, namely\n\\begin{equation}\n2F(R_0)-R_0F_R(R_0)=0\\,,\\label{trace}\n\\end{equation}\nwhich leads in our case,\n\\begin{equation}\nR_0=\\left(\\frac{1}{\\alpha(n-2)}\\right)^{\\frac{1}{n-1}}\\,.\\label{RdS}\n\\end{equation}\nThe model with $n=2$ does not possess an exact de Sitter solution, but if \n$n>2$, we can realize it. Here, we remember that $\\alpha>0$.\nIn the Einstein frame this solution corresponds to\n\\begin{equation}\n\\sigma_{max}=-\\sqrt{\\frac{3}{2\\kappa^2}}\\log\\left[\\frac{2n-2}{n-2}\\right]\\,,\n\\end{equation}\nwhere we have used (\\ref{sigma}). We note that $\\sigma<0$ only if $n>2$ (we are \nassuming $n>1$); otherwise, if $n<2$, this expression becomes imaginary and \nmeaningless.\nThis solution corresponds to a maximum of the potential, because of\n\\begin{equation}\nV'(\\sigma_{max})=0\\,,\\quad \nV''(\\sigma_{max})=\\frac{1}{6}\\left(\\frac{1}{\\alpha\\,n}\\right)^{\\frac{1}{n-1}}\n\\frac{n-2}{n-1}\\left(\\frac{n}{n-2}\\right)^{\\frac{(2-n)}{(n-1)}}>0\\,.\n\\end{equation}\nWhen $n>2$, the potential in Einstein frame has a maximum. The scalar field \nproduces the de Sitter solution, but the field does not evolve with the time \n($\\dot\\sigma=0$) and we do not have a natural exit form inflation. However, by \nmaking use of the perturbative theory, we may study the stability of the model \nin the Jordan frame.\n\nBy perturbating Eq.~(\\ref{eq}) as $R\\rightarrow R+\\delta R$, $|\\delta R|\\ll 1$ \naround the de Sitter solution, we get\n\\begin{equation}\n\\left(\\Box-m^2\\right)\\delta R\\simeq 0\\,,\n\\quad\nm^2=\\frac{1}{3}\\left(\\frac{F'(R)}{F''(R)}-R\n\\right)\\,.\n\\end{equation}\nHere, $m^2$ is the effective mass of the scalaron, namely the new degree of \nfreedom\nintroduced by the modified gravity through $F_R(R)$, which is proportional to \nthe opposite of the inflaton, namely $-\\sigma$. As a consequence, $m^2$ is \nproportional to the opposite of the scalar potential of the inflaton and when \n$m^2<0$ the solution is unstable. In our case, we get on the de Sitter \nsolution,\n\\begin{equation}\nm^2=-(n-2)^{\\frac{n-2}{n-1}}\\left(\\frac{1}{\\alpha\\,n}\\right)^{\\frac{1}{n-1}}<0\\,,\\quad \n22\\,,\n\\end{equation}\nsince it is easy to demonstrate that $x<0$ if $0<16(n-2)\/n$, that is always \ntrue when $20$, the Ricci scalar decreases with the red shift and the physics of \nStandard Model can emerge.\nIn Ref.~\\cite{nostrum} numerical calculations have been executed in different \ninflationary models\nprovided by the specific $R^n$-term with $22$, may represent a valid inflationary scenario. \nIn this case, the curvature is bounded at a specific value at the beginning of \ninflation, and the de Sitter space-time is an exact solution of the model, \nwhich results to be highly unstable. From (\\ref{delta}) we see how extremely \nsmall perturbations in hot universe give the correct $N$-foldings number and \nthe exit from inflation.\nFinally, it is important to note that positive energy density of perturbations \nbrings the curvature to decrease during the early-time acceleration making the \nmodel consistent with the observable evolution of our universe.\n\n\n\\section{Conclusions}\n\n\nThe attention that recently has been paid to modified theories of gravity is \ncaused by the idea of unified description of early-time and late-time cosmic \nacceleration~\\cite{Review-Nojiri-Odintsov}. Moreover,\nthe gravitational action of such a kind of theories may describe quantum \neffects in hot universe scenario, and the last cosmological data seems to be in \nfavour of quadratic corrections to General Relativity during this phase.\n\nIn this paper, we have investigated some feautures of $F(R)$-modified gravity \nmodels for inflationary cosmology, by performing our analysis in the Jordan and \nin the Einstein framework.\n\nAt first, we have studied\ninflation for the class of scalar potentials of the type \n$V(\\sigma)\\sim\\exp[n\\sigma]$, $n$ being a general parameter, in Einstein frame. \nAs a matter of fact, for such a kind of models is possible to reconstruct the \n$F(R)$-gravity theories which correspond to the given potentials.\nSince\nviable inflation must be consistent with the last Planck data, the potentials \nhave been carefully alalyzed, by finding the conditions on the parameters which \nmake possible the early-time acceleration according with $N$-foldings, spectral\nindex and tensor-to-scalar ratio coming from observations.\nWe have derived the form of the $F(R)$-models at the large and small curvature \nlimits, demonstrating that these\nmodels in the Jordan frame correspond to corrections to Einstein gravity which \nemerge only at mass scales larger than the Planck mass.\nThe investigated potentials can be classified in two classes. In the first one, \nthe scalar potential behaves as \n$V(\\sigma)\\sim\\exp\\left[n\\sqrt{2\\kappa^2\/3}\\sigma\\right]\\,,n>0$, and produces \nacceleration with a sufficient amount of inflation only for $n$ very close to \nzero, namely the $F(R)$-model must be closed to the Starobinsky one. In the \nsecond class, $V(\\sigma)\\sim\\alpha \n-\\gamma\\exp\\left[-n\\sqrt{2\\kappa^2\/3}\\sigma\\right]$, $n>0$, and the quasi de \nSitter solution emerges during viable inflation when $n>1\/3$. This kind of \nscalar potentials is originated from large curvature corrections of the type \n$F(R)\\sim c_1 R^2+c_2 R^\\zeta$, with $1<\\zeta<5\/3$ when $1\/31$.\nIt is interesting to note that $R^2$-term which induces the early-time \ninflation is also responsible for removal of finite-time future \nsingularity in $F(R)$-gravity unifying the inflation with dark energy~\\cite{OdSing}.\n\nIn the second part of the paper, we have studied in detail the specific class \nof models $F(R)=R+(R+R_0)^n$. The analysis in the\nEinstein frame reveals that $n$ must be very close to two in order to realize a \nviable inflation for large and negative values of the scalar field,\nbut other possibilities are allowed\nby starting from intermediate values of the field.\nTo be specific, when $n>2$, the de Sitter solution emerges, but in order to \nstudy the exit from inflation, since in this case the de Sitter space-time is \nan exact solution and the field does not move and exit from inflation in a \nnatural way, is necessary to analyze the theory in the Jordan frame, where \nperturbations make possible an early time acceleration with a sufficient amount \nof inflation. Moreover, we can explicitly demonstrate\nthat the curvature decreases making the model consistent with the historical \nevolution of our universe.\n\n\n\n\\section*{Acknowledgments}\nWe would like to thank M. Sami for useful discussions and valuable suggestions.\nThe work by SDO has been supported in part by MINECO (Spain), project \nFIS2010-15640, by AGAUR (Generalitat de Catalunya), contract 2009SGR-994 and by \nproject 2.1839.2011 of MES (Russia).\n\n\n\n\n\\section*{Appendix A}\nLet us consider the equation\n\\begin{equation}\nF_R(R)-\\frac{1}{2}-\\frac{F_R(R)^{1+n}}{2} =-\\frac{R}{4\\gamma(n+2)}\\,,\n\\label{EqA0}\\end{equation}\nwith $n\\,,\\gamma>0$ (it corresponds to (\\ref{pippo}) with (\\ref{alphacond})).\nWe want to find solutions of the above equation as a function of $R$.\nA simple possibility is looking for regular solutions as a power series of $R$, \nnamely\nsetting\n\\begin{equation*}\nF_R(R)=1+c_1R+c_2R^2+c_3R^3+...\\,,\n\\end{equation*}\nwhere the constants $c_{1,2,3...}$ have to be determined. At first, we are \ninterested in solutions where $|c_{1,2,3...}R^{1,2,3...}|\\ll 1$, namely we \nanalyze the limit at small curvature. Thus,\nplugging this Ansatz in the above equation one has\n\\begin{equation*}\n(1+c_1R+c_2R^2+c_3R^3+...)-\\frac{1}{2}-\\frac{1}{2}\\left(1+c_1R+c_2R^2+ \nc_3R^3+... \\right)^{1+n}=\n-\\frac{R}{4\\gamma(n+2)}\\,.\n\\end{equation*}\nPutting\n\\begin{equation*}\nX=+c_1R+c_2R^2+ c_3R^3+... \\,,\n\\end{equation*}\nand assuming $|X|\\ll 1$, the following expansion holds\n\\begin{equation*}\n(1+X)^{1+n}=1+(1+n)X+\\frac{(1+n)n}{2}X^2+\\frac{n(n+1)(n-1)}{3!}X^3+...\\,.\n\\end{equation*}\nAs a consequence, one arrives at the recursive relations\n\\begin{equation*}\nc_1=\\frac{1}{2\\gamma(n-1)(n+2)}\\,,\\quad\n c_2=-\\frac{(1+n)n\\,c_1^2}{2(n-1)}\\,,\\quad\n c_3=-\\frac{(1+n)c_1^3}{6}\\,, ...\\,,\n\\end{equation*}\nfrom which it follows\n\\begin{equation*}\nc_1 R\\propto\\frac{R}{\\gamma}\\,,\\quad\n c_2 R^2\\propto\\frac{R}{\\gamma^2}\\,,\\quad\n c_3 R^3\\propto\\frac{R^3}{\\gamma^3}\\,,...\n\\end{equation*}\n This approximation is valid when $R\\ll\\gamma$.\nOn the other side, when $R\\gg\\gamma$ we can check for the solutions in the \nform\n\\begin{equation*}\nF_R(R)=c_0\\left(\\frac{R}{\\gamma}\\right)^\\zeta\\,,\\quad\\zeta<1\\,.\n\\end{equation*}\nIn such a case from (\\ref{EqA0}) we get\n\\begin{equation*}\nc_0^{1+n}\\left(\\frac{R}{\\gamma}\\right)^{\\zeta(1+n)}\\simeq\\frac{1}{4(n+2)}\\left(\\frac{R}{\\gamma}\\right)\\,,\n\\end{equation*}\nand as a consequence\n\\begin{equation*}\nc_0=\\left(\\frac{1}{4(n+2)}\\right)^{\\frac{1}{1+n}}\\,,\\quad\\zeta=\\frac{1}{1+n}<1\\,.\n\\end{equation*}\nNow let us consider the following equation (it corresponds to \\ref{pippo2} \nwith \\ref{alphaalpha}),\n\\begin{equation}\nF_R(R)-F_R(R)^{1-n} =\\frac{R}{2\\gamma(2-n)}\\,,\n\\label{EqA1}\\end{equation}\nwhere $\\gamma>0$ and $0\\alpha_2\\,.\n\\end{equation*}\nSince $0 0$.\nChoose an open set $U \\S X$ such that $x_0 \\in U$ and\nfor all $x \\in U$ we have\n\\[\n\\| a (x) - a (x_0) \\| < \\frac{\\varepsilon}{2}\n\\qquad {\\mbox{and}} \\qquad\n\\| \\alpha_x (a (x_0)) - \\alpha_{x_0} (a (x_0)) \\| < \\frac{\\varepsilon}{2}.\n\\]\nThen, using $\\| \\alpha_x \\| = 1$ for all $x \\in X$,\none sees that $x \\in U$\nimplies $\\| \\alpha_x (a (x)) - \\alpha_{x_0} (a (x_0)) \\| < \\varepsilon$.\n\\end{proof}\n\nFor any group~$G$,\nthere are also algebraic conditions\nrelating the automorphisms $\\alpha_{g, x}$,\ncoming from the requirement that $g \\mapsto \\alpha_g$\nbe a group homomorphism.\nIf $G = {\\mathbb{Z}}$,\nthen we really need only the function $x \\mapsto \\alpha_{1, x}$.\nFor reference,\nwe give the relevant statement as a lemma.\n\n\\begin{lem}\\label{L_4Y12_GetAuto}\nLet $X$ be a locally compact Hausdorff space,\nlet $h \\colon X \\to X$ be a homeomorphism,\nlet $D$ be a C*-algebra.\nThen there is a one to one correspondence between actions of ${\\mathbb{Z}}$ on\n$C_0 (X, D)$\nthat lie over $h$ and\ncontinuous functions from $X$ to ${\\mathrm{Aut}} (D)$,\ngiven as follows.\n\nFor any function $x \\mapsto \\alpha_{x}$\nfrom $X$ to ${\\mathrm{Aut}} (D)$\nsuch that $x \\mapsto \\alpha_{x} (d)$ is continuous{} for all $d \\in D$,\nthere is an automorphism $\\alpha \\in {\\mathrm{Aut}} ( C_0 (X, D))$\ngiven by $\\alpha (a) (x) = \\alpha_x (a (h^{-1} (x))$\nfor all $a \\in C_0 (X, D)$ and $x \\in X$,\nand this automorphism lies over~$h$.\n\nConversely,\nif $\\alpha \\in {\\mathrm{Aut}} ( C_0 (X, D))$ lies over~$h$,\nthen there is a function $x \\mapsto \\alpha_{x}$\nfrom $X$ to ${\\mathrm{Aut}} (D)$\nsuch that $x \\mapsto \\alpha_{x} (d)$ is continuous{} for all $d \\in D$\nand such that $\\alpha (a) (x) = \\alpha_x (a (h^{-1} (x))$\nfor all $a \\in C_0 (X, D)$ and $x \\in X$.\n\\end{lem}\n\n\\begin{proof}\nUsing Lemma~\\ref{L_9Y30_MixCont}, this is immediate.\n\\end{proof}\n\nThere is a conflict in the notation in Lemma~\\ref{L_4Y12_GetAuto}:\nif $n \\in {\\mathbb{Z}}$\nthen $\\alpha_n$\nis one of the automorphisms in the action on $C_0 (X, D)$\n(namely $\\alpha_n$),\nwhile if $x \\in X$\nthen $\\alpha_x \\in {\\mathrm{Aut}} (D)$.\nWe use this notation anyway to avoid having more letters.\nTo distinguish the two uses,\ntake $\\alpha$ to be the action\nand write $x \\mapsto \\alpha_x$\nwhen the function from $X$ to ${\\mathrm{Aut}} (D)$ is intended.\n\nIf $D$ is prime,\nthen every action on $C_0 (X, D)$\nlies over an action of $G$ on~$X$.\n\n\\begin{lem}\\label{N_4Y12_IfDSimple}\nLet $X$ be a locally compact Hausdorff space,\nlet $D$ be a prime C*-algebra,\nlet $G$ be a topological group,\nand let $\\alpha \\colon G \\to {\\mathrm{Aut}} (C_0 (X, D))$\nbe an action of $G$ on $C_0 (X, D)$.\nThen there exists an action of $G$ on~$X$\nsuch that $\\alpha$ lies over this action.\n\\end{lem}\n\n\\begin{proof}\nThe action of $G$ on~$X$\nis obtained from the identification\n$X \\cong {\\mathrm{Prim}} ( C_0 (X, D) )$.\n\\end{proof}\n\n\\begin{prp}\\label{P_4Y15_CPSimple}\nLet $G$ be a discrete group,\nlet $X$ be a compact space,\nand suppose $G$ acts on $X$ in such a way that\nthe action is minimal and\nfor every finite set $S \\subset G \\setminus \\{ 1 \\}$,\nthe set\n\\[\n\\big\\{ x \\in X \\colon {\\mbox{$g x \\neq x$ for all $g \\in S$}} \\big\\}\n\\]\nis dense in~$X$.\nLet $D$ be a simple unital C*-algebra,\nand let $\\alpha \\colon G \\to {\\mathrm{Aut}} (C (X, D))$\nbe an action of $G$ on $C (X, D)$\nwhich lies over the given action of $G$ on $X$\n(in the sense of Definition~\\ref{D_4Y12_OverX}).\nThen $C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nis simple.\n\\end{prp}\n\n\\begin{proof}\nFor any C*-algebra~$A$,\nlet ${\\widehat{A}}$ be the space of unitary equivalence classes\nof irreducible representations of~$A$,\nwith the hull-kernel topology.\nSince the primitive ideals\nof $C (X, D)$ are exactly the kernels of the point evaluations,\nthere is an obvious map $q \\colon \\CXDHat{X}{D} \\to X$,\nand the open sets in $\\CXDHat{X}{D}$ are exactly\nthe sets $q^{-1} (U)$ for open sets $U \\subset X$.\nIt is now immediate that for every finite set\n$S \\subset G \\setminus \\{ 1 \\}$,\nthe set\n\\[\n\\big\\{ x \\in \\CXDHat{X}{D} \\colon\n {\\mbox{$g x \\neq x$ for all $g \\in S$}} \\big\\}\n\\]\nis dense in $\\CXDHat{X}{D}$.\nThat is, the action of $G$ on $\\CXDHat{X}{D}$\nis topologically free in the sense of Definition~1 of~\\cite{AS}.\n\nLet $J \\subset C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nbe a nonzero ideal.\nLet\n\\[\n\\pi \\colon C^* \\big( G, \\, C (X, D), \\, \\alpha \\big)\n \\to C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)\n\\]\nbe the quotient map.\nTheorem~1 of~\\cite{AS}\nimplies that $\\pi^{-1} (J)$ has nonzero intersection\nwith the canonical copy of $C (X, D)$\nin $C^* \\big( G, \\, C (X, D), \\, \\alpha \\big)$.\nTherefore $J$ has nonzero intersection\nwith the canonical copy of $C (X, D)$\nin $C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$.\nSince $D$ is simple,\nthis intersection has the form $C_0 (U, D)$\nfor some nonempty open set $U \\subset X$.\n\nSince the action of $G$ on $X$ is minimal and $X$ is compact,\nthere exist $n \\in {\\mathbb{Z}}_{> 0}$ and $g_1, g_2, \\ldots, g_n \\in {\\mathbb{Z}}_{> 0}$\nsuch that the sets\n$g_1^{-1} U, \\, g_2^{-1} U, \\, \\ldots, g_n^{-1} U$ cover~$X$.\nChoose\n\\[\nf_1, f_2, \\ldots, f_n\n \\in C (X)\n \\subset C (X, D)\n \\subset C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)\n\\]\nsuch that ${\\mathrm{supp}} (f_k) \\subset g_k^{-1} U$ for $k = 1, 2, \\ldots, n$\nand $\\sum_{k = 1}^n f_k = 1$.\nFor $k = 1, 2, \\ldots, n$, the functions\n$\\alpha_{g_k} (f_k)$\nare in $C_0 (U) \\subset C_0 (U, D) \\subset J$,\nso\n\\[\n1 = \\sum_{k = 1}^n f_k\n = \\sum_{k = 1}^n u_{g_k}^* \\alpha_{g_k} (f_k) u_{g_k}\n \\in J.\n\\]\nSo $J = C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$.\n\\end{proof}\n\n\\begin{prp}\\label{P_5418_CPStFin}\nAssume the hypotheses of Proposition~\\ref{P_4Y15_CPSimple},\nand in addition assume that\n$G$ is amenable\nand $D$ has a tracial state.\nThen $C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nhas a tracial state\nand is stably finite.\n\\end{prp}\n\n\\begin{proof}\nSince $C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nis simple by Proposition~\\ref{P_4Y15_CPSimple},\nit suffices to show that\n$C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nhas a tracial state.\nWe know that the tracial state space\n${\\operatorname{T}} (C (X, D))$ is nonempty,\nsince one can compose a tracial state on~$D$\nwith a point evaluation\n$C (X, D) \\to D$.\nSince $G$ is amenable,\ncombining Theorem 2.2.1 and 3.3.1 of \\cite{Gr}\nshows that $C (X, D)$ has a $G$-invariant tracial state~$\\tau$.\nStandard methods show that the composition of $\\tau$\nwith the conditional expectation\nfrom $C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nto $C (X, D)$\nis a tracial state on\n$C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$.\n\\end{proof}\n\nThe following condition is a technical hypothesis\nwhich we need for the proof of the main large subalgebra result\n(Theorem~\\ref{T_5418_AYLg}).\n\n\\begin{dfn}\\label{D_4Y12_PsPer}\nLet $D$ be a C*-algebra,\nand let $S \\subset {\\mathrm{Aut}} (D)$.\nWe say that $S$ is {\\emph{pseudoperiodic}}\nif for every $a \\in D_{+} \\setminus \\{ 0 \\}$\nthere is $b \\in D_{+} \\setminus \\{ 0 \\}$\nsuch that for every $\\alpha \\in S \\cup \\{ {\\mathrm{id}}_D \\}$\nwe have $b \\precsim \\alpha (a)$.\n\\end{dfn}\n\nThe interpretation of pseudoperiodicity is roughly as follows.\nSuppose $S \\subset {\\mathrm{Aut}} (D)$ is pseudoperidic.\nThen there is no sequence $(\\alpha_n)_{n \\in {\\mathbb{Z}}_{> 0}}$\nin $S$ for which there is a nonzero element\n$\\eta \\in W (D)$\nsuch that the sequence $(\\alpha_{n} (\\eta))_{n \\in {\\mathbb{Z}}_{> 0}}$\nbecomes arbitrarily small in $W (D)$\nin a heuristic sense.\n\nWe give some conditions which imply pseudoperiodicity.\n\n\\begin{lem}\\label{L_4Y12_ApproxInn}\nLet $D$ be a unital C*-algebra.\nThen the set of approximately inner automorphisms of $D$\nis pseudoperiodic in the sense of Definition~\\ref{D_4Y12_PsPer}.\n\\end{lem}\n\n\\begin{proof}\nLet $a \\in D_{+} \\setminus \\{ 0 \\}$.\nIt suffices to prove that\n$a \\precsim \\alpha (a)$ for every approximately inner automorphism\n$\\alpha \\in {\\mathrm{Aut}} (D)$.\nTo see this, let $\\varepsilon > 0$ be arbitrary,\nand use approximate innerness\nto chose a unitary $u \\in D$ such that\n$\\| u \\alpha (a) u^* - a \\| < \\varepsilon$.\n\\end{proof}\n\n\\begin{lem}\\label{Lem1_190201D}\nLet $D$ be a simple C*-algebra.\nLet $S \\S {\\mathrm{Aut}} (D)$\nbe a subset which is compact in the topology\nof pointwise convergence in the norm on~$D$.\nThen $S$ is pseudoperiodic\nin the sense of Definition~\\ref{D_4Y12_PsPer}.\n\\end{lem}\n\n\\begin{proof}\nLet $a \\in D_{+} \\setminus \\{ 0 \\}$.\nWithout loss of generality{} $\\| a \\| = 1$.\nFor $\\beta \\in S \\cup \\{ {\\mathrm{id}}_D \\}$ set\n\\[\nU_{\\beta}\n = \\left\\{ \\alpha \\in S \\cup \\{ {\\mathrm{id}}_D \\} \\colon\n \\| \\alpha (a) - \\beta (a) \\| < \\frac{1}{2} \\right\\}.\n\\]\nBy compactness,\nthere are $\\beta_1, \\beta_2, \\ldots, \\beta_n \\in S \\cup \\{ {\\mathrm{id}}_D \\}$\nsuch that $U_{\\beta_1}, U_{\\beta_2}, \\ldots, U_{\\beta_n}$\ncover $S \\cup \\{ {\\mathrm{id}}_D \\}$.\nSince $\\bigl( \\beta (a) - \\frac{1}{2} \\bigr)_{+} \\neq 0$\nfor all $\\beta \\in {\\mathrm{Aut}} (D)$,\nby Lemma~2.6 of~\\cite{PhLg} there is $b \\in D_{+} \\setminus \\{ 0 \\}$\nsuch that $b \\precsim \\bigl( \\beta_k (a) - \\frac{1}{2} \\bigr)_{+}$\nfor $j = 1, 2, \\ldots, n$.\nLet $\\alpha \\in S \\cup \\{ {\\mathrm{id}}_D \\}$.\nChoose $k \\in \\{ 1, 2, \\ldots, n \\}$ such that $\\alpha \\in U_{\\beta_k}$.\nThen $\\| \\beta_k (a) - \\alpha (a) \\| < \\frac{1}{2}$,\nso\n$b \\precsim \\bigl( \\beta_k (a) - \\frac{1}{2} \\bigr)_{+}\n \\precsim \\alpha (a)$.\n\\end{proof}\n\nThe following result will not be used,\nsince large subalgebras are not used in our proofs\nwhen $D$ is purely infinite.\nIt is included as a further example of pseudoperiodicity.\n\n\\begin{lem}\\label{L_9212_PISCA}\nLet $D$ be a purely infinite simple C*-algebra.\nThen ${\\mathrm{Aut}} (D)$\nis pseudoperiodic in the sense of Definition~\\ref{D_4Y12_PsPer}.\n\\end{lem}\n\n\\begin{proof}\nLet $a \\in D_{+} \\setminus \\{ 0 \\}$.\nThen $a \\precsim b$ for all $b \\in D_{+} \\setminus \\{ 0 \\}$.\nIn particular,\n$a \\precsim \\alpha (a)$ for all $\\alpha \\in {\\mathrm{Aut}} (D)$.\n\\end{proof}\n\n\\begin{lem}\\label{L_4Y12_TrPrsv}\nLet $D$ be a simple unital C*-algebra{}\nwhich has strict comparison of positive elements.\nThen ${\\mathrm{Aut}} (D)$\nis pseudoperiodic in the sense of Definition~\\ref{D_4Y12_PsPer}.\n\\end{lem}\n\n\\begin{proof}\nIf $D$ is finite dimensional,\nthe conclusion is immediate.\nOtherwise,\nlet $a \\in D_{+} \\setminus \\{ 0 \\}$.\nWithout loss of generality{} $\\| a \\| \\leq 1$.\nThen $\\tau (a) \\leq d_{\\tau} (a)$\nfor all $\\tau \\in {\\operatorname{QT}} (D)$.\nMoreover,\nsince ${\\operatorname{QT}} (D)$ is compact and $D$ is simple,\nthe number\n$\\delta = \\inf_{\\tau \\in {\\operatorname{QT}} (D)} \\tau (a)$\nsatisfies $\\delta > 0$.\nUse Corollary~2.5 of~\\cite{PhLg}\nto find\n$b \\in D_{+} \\setminus \\{ 0 \\}$\nsuch that\n$d_{\\tau} (\\langle b \\rangle ) < \\delta$\nfor all $\\tau \\in {\\operatorname{QT}} (D)$.\nThen\nfor every $\\tau \\in {\\operatorname{QT}} (D)$,\nusing $\\tau \\circ \\alpha \\in {\\operatorname{QT}} (D)$\nat the second step,\nwe have\n\\[\nd_{\\tau} (b)\n < \\delta\n \\leq (\\tau \\circ \\alpha) (a)\n \\leq d_{\\tau} ( \\alpha (a) ).\n\\]\nThe strict comparison hypothesis therefore implies\nthat $b \\precsim_A \\alpha (a)$.\n\\end{proof}\n\n\\begin{dfn}\\label{Def2-181221D}\nLet $A$ be a C*-algebra.\nWe say that the\n{\\emph{order on projections over $A$\nis determined by quasitraces}}\nif whenever $p,q \\in M_{\\infty} (A)$ are projections such that\n$\\tau (p) < \\tau(q) $\nfor all $\\tau \\in \\operatorname{QT} (A)$, then $p \\precsim q$.\n\\end{dfn}\n\n\\begin{lem}\\label{L_5417_SP}\nLet $D$ be a simple unital C*-algebra{}\nwith Property~(SP)\nand such that\nthe order on projections over~$A$ is determined by quasitraces.\nThen ${\\mathrm{Aut}} (D)$\nis pseudoperiodic in the sense of Definition~\\ref{D_4Y12_PsPer}.\n\\end{lem}\n\n\\begin{proof}\nIf $D$ is finite dimensional,\nthe conclusion is immediate.\nOtherwise,\nlet $a \\in D_{+} \\setminus \\{ 0 \\}$.\nChoose a nonzero projection{} $p \\in {\\overline{a D a}}$.\nSince ${\\operatorname{QT}} (D)$ is compact and $D$ is simple,\nthe number\n$\\delta = \\inf_{\\tau \\in {\\operatorname{QT}} (D)} \\tau (p)$\nsatisfies $\\delta > 0$.\nChoose $n \\in {\\mathbb{Z}}_{> 0}$\nsuch that $\\frac{1}{n} < \\delta$.\nUse Lemma~2.3 of~\\cite{PhLg}\nto choose a unitary $u \\in A$\nand a nonzero positive element $b \\in A$\nsuch that the elements\n\\[\nb, \\, u b u^{-1}, \\, u^2 b u^{-2}, \\, \\ldots, \\, u^n b u^{-n}\n\\]\nare pairwise orthogonal.\nChoose a nonzero projection{} $q \\in {\\overline{b D b}}$.\nThen\nfor every $\\tau \\in {\\operatorname{QT}} (D)$,\nusing $\\tau \\circ \\alpha \\in {\\operatorname{QT}} (D)$\nat the third step,\nwe have\n\\[\n\\tau (q)\n \\leq \\frac{1}{n + 1}\n < \\delta\n \\leq (\\tau \\circ \\alpha) (p).\n\\]\nThe strict comparison hypothesis therefore implies\nthat $q$ is Murray-von Neumann equivalent{} to a subprojection of~$\\alpha (p)$.\nIt follows that $q \\precsim_A \\alpha (a)$.\n\\end{proof}\n\n\nThe following definition is intended only for convenience\nin this paper.\n(The condition occurs several times as a hypothesis,\nand is awkward to state.)\n\n\\begin{dfn}\\label{D_4Y16_PsPrGen}\nLet $X$ be a locally compact Hausdorff space,\nlet $G$ be a topological group,\nand let $D$ be a C*-algebra.\nLet $\\alpha \\colon G \\to {\\mathrm{Aut}} (C_0 (X, D))$\nbe an action of $G$ on $C_0 (X, D)$\nwhich lies over an action of $G$ on~$X$,\nand let $(g, x) \\mapsto \\alpha_{g, x} \\in {\\mathrm{Aut}} (D)$\nbe as in Definition~\\ref{D_4Y12_OverX}.\nWe say that $\\alpha$ is\n{\\emph{pseudoperiodically generated}}\nif\n\\[\n\\big\\{ \\alpha_{g, x} \\colon {\\mbox{$g \\in G$ and $x \\in X$}} \\big\\}\n\\]\nis pseudoperiodic in the sense of Definition~\\ref{D_4Y12_PsPer}.\n\\end{dfn}\n\n\\begin{lem}\\label{L_4Y16_WhenPPG}\nLet $X$ be a compact metric space, let $h \\colon X \\to X$ be a homeomorphism,\nlet $D$ be a simple unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nAs in Lemma~\\ref{L_4Y12_GetAuto},\nlet $( \\alpha_{x} )_{x \\in X}$\nbe the family in ${\\mathrm{Aut}} (D)$\nsuch that $\\alpha (a) (x) = \\alpha_x (a (h^{-1} (x))$\nfor all $a \\in C_0 (X, D)$ and $x \\in X$.\nSuppose that the subgroup $H$ of ${\\mathrm{Aut}} (D)$ generated by\n$\\{ \\alpha_x \\colon x \\in X \\}$ is pseudoperiodic.\nThen the action of ${\\mathbb{Z}}$ generated by $\\alpha$\nis pseudoperiodically generated.\n\\end{lem}\n\n\\begin{proof}\nOne checks that if $(n, x) \\mapsto \\alpha_{n, x} \\in {\\mathrm{Aut}} (D)$\nis determined\n(following the notation of Definition~\\ref{D_4Y12_OverX})\nby $\\alpha^n (a) (x) = \\alpha_{n, x} (a (h^{-n} (x))$,\nthen $\\alpha_{n, x} \\in H$ for all $n \\in {\\mathbb{Z}}$ and $x \\in X$.\n\\end{proof}\n\n\\section{The orbit breaking subalgebra for a nonempty set meeting\neach orbit at most once}\\label{Sec_OrbBreak}\n\n\\indent\nLet $h \\colon X \\to X$ be a homeomorphism{}\nof a compact Hausdorff space~$X$,\nand let $D$ be a simple unital C*-algebra.\nFor $Y \\subset X$ closed,\nfollowing Putnam~\\cite{Pt1},\nin Definition~7.3 of~\\cite{PhLg}\nwe defined the $Y$-orbit breaking subalgebra\n$C^* ({\\mathbb{Z}}, X, h)_Y \\subset C^* ({\\mathbb{Z}}, X, h)$.\nHere,\nfor an automorphism $\\alpha \\in {\\mathrm{Aut}} ( C (X, D) )$\nwhich lies over~$h$\nwe define $C^* ({\\mathbb{Z}}, C (X, D), \\alpha)_Y$.\nWe prove that if $X$ is infinite,\n$h$~is minimal,\n$Y$ intersects each orbit at most once,\nand an additional technical condition is satisfied\n(namely, that the action of ${\\mathbb{Z}}$ generated by~$\\alpha$\nis pseudoperiodically generated),\nthen $C^* ({\\mathbb{Z}}, C (X, D), \\alpha)_Y$\nis a large subalgebra of $C^* ({\\mathbb{Z}}, C (X, D), \\alpha)$\nof crossed product type,\nin the sense of Definition~4.9 of~\\cite{PhLg}.\nThis is a generalization\nof Theorem~7.10 of~\\cite{PhLg}.\n\n\n\\begin{ntn}\\label{N_4Y12_CPN}\nLet $G$ be a discrete group,\nlet $A$ be a C*-algebra,\nand let $\\alpha \\colon G \\to {\\mathrm{Aut}} (A)$ be an action of $G$ on~$A$.\nWe identify $A$ with a subalgebra of $C^*_{\\mathrm{r}} (G, A, \\alpha)$\nin the standard way.\nWe let $u_g \\in M (C^*_{\\mathrm{r}} (G, A, \\alpha))$\nbe the standard unitary\ncorresponding to $g \\in G$.\nWhen $G = {\\mathbb{Z}}$, we write just $u$ for the unitary~$u_1$\ncorresponding to the generator $1 \\in {\\mathbb{Z}}$.\nWe let $A [G]$ denote the dense *-subalgebra\nof $C^*_{\\mathrm{r}} (G, A, \\alpha)$\nconsisting of sums $\\sum_{g \\in S} a_g u_g$\nwith $S \\subset G$ finite and $a_g \\in A$ for $g \\in S$.\nWe may always assume $1 \\in S$.\nWe let $E_{\\alpha} \\colon C^*_{\\mathrm{r}} (G, A, \\alpha) \\to A$\ndenote the standard conditional expectation,\ndefined on $A [G]$ by\n$E_{\\alpha} \\left( \\sum_{g \\in S} a_g u_g \\right) = a_1$.\nWhen $\\alpha$ is understood, we just write~$E$.\n\nWhen $G$ acts on a compact Hausdorff space~$X$,\nwe use obvious analogs of this notation for $C^*_{\\mathrm{r}} (G, X)$,\nwith the action as in Notation~\\ref{N_5418_ActCX}.\nIn particular,\nif $G = {\\mathbb{Z}}$ and the action of generated by a homeomorphism{}\n$h \\colon X \\to X$,\nwe have $u f u^* = f \\circ h^{-1}$.\n\\end{ntn}\n\n\\begin{ntn}\\label{N_4Y12_C0UD}\nFor a locally compact Hausdorff space~$X$\nand a C*-algebra~$D$,\nwe identify $C_0 (X, D) = C_0 (X) \\otimes D$\nin the standard way.\nFor an open subset $U \\subset X$,\nwe use the abbreviation\n\\[\nC_0 (U, D)\n = \\big\\{ a \\in C_0 (X, D) \\colon\n {\\mbox{$a (x) = 0$ for all $x \\in X \\setminus U$}} \\big\\}\n \\subset C_0 (X, D).\n\\]\nThis subalgebra is of course canonically isomorphic to\nthe usual algebra $C_0 (U, D)$ when $U$ is considered\nas a locally compact Hausdorff space{} in its own right.\n\\end{ntn}\n\nIn particular,\nif $Y \\subset X$ is closed, then\n\\begin{equation}\\label{Eq_4Y12_C0XYD}\nC_0 (X \\setminus Y, \\, D)\n = \\big\\{ a \\in C_0 (X, D) \\colon\n {\\mbox{$a (x) = 0$ for all $x \\in Y$}} \\big\\}.\n\\end{equation}\n\nThe following definition is the analog of\nDefinition~7.3 of~\\cite{PhLg}.\n\n\\begin{dfn}\\label{D_4Y12_OrbSubalg}\nLet $X$ be a locally compact Hausdorff space,\nlet $h \\colon X \\to X$ be a homeomorphism,\nlet $D$ be a C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} ( C_0 (X, D))$\nbe an automorphism which lies over~$h$.\nLet $Y \\subset X$ be a nonempty closed subset,\nand, following~(\\ref{Eq_4Y12_C0XYD}), define\n\\[\nC^* \\big( {\\mathbb{Z}}, \\, C_0 (X, D), \\, \\alpha \\big)_Y\n = C^* \\big( C_0 (X, D), \\, C_0 (X \\setminus Y, \\, D) u \\big)\n \\subset C^* \\big( {\\mathbb{Z}}, \\, C_0 (X, D), \\, \\alpha \\big).\n\\]\nWe call it the {\\emph{$Y$-orbit breaking subalgebra}}\nof $C^* \\big( {\\mathbb{Z}}, C_0 (X, D), \\alpha \\big)$.\n\\end{dfn}\n\nWe show that if $Y$ intersects each orbit of~$h$ at most once,\nand the action of ${\\mathbb{Z}}$ generated by~$\\alpha$\nis pseudoperiodically generated,\nthen $C^* ({\\mathbb{Z}}, C_0 (X, D), \\alpha)_Y$\nis a large subalgebra of $C^* \\big( {\\mathbb{Z}}, \\, C_0 (X, D), \\, \\alpha \\big)$\nof crossed product type.\n\nThe following lemma is the analog of Proposition~7.5 of~\\cite{PhLg}.\n\n\\begin{lem}\\label{L_4X21_AYStruct}\nLet $X$ be a compact Hausdorff space, let $h \\colon X \\to X$ be a homeomorphism,\nlet $D$ be a unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nLet\n\\[\nu \\in C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)\n\\qquad {\\mbox{and}} \\qquad\nE_{\\alpha} \\colon C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr) \\to C (X, D)\n\\]\nbe as in Notation~\\ref{N_4Y12_CPN}.\nLet $Y \\subset X$ be a nonempty closed subset.\nFor $n \\in {\\mathbb{Z}}$, set\n\\[\nY_n = \\begin{cases}\n \\bigcup_{j = 0}^{n - 1} h^j (Y) & \\hspace{3em} n > 0\n \\\\\n \\varnothing & \\hspace{3em} n = 0\n \\\\\n \\bigcup_{j = 1}^{- n} h^{-j} (Y) & \\hspace{3em} n < 0.\n\\end{cases}\n\\]\nThen\n\\begin{align}\\label{Eq:2816CharOB}\n& C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y\n\\\\\n& = \\big\\{ a \\in C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr) \\colon\n {\\mbox{$E_{\\alpha} (a u^{-n}) \\in C_0 (X \\setminus Y_n, \\, D)$\n for all $n \\in {\\mathbb{Z}}$}} \\big\\}\n \\notag\n\\end{align}\nand\n\\begin{equation}\\label{Eq:2816Dense}\n{\\overline{C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y \\cap\nC (X, D) [{\\mathbb{Z}}]}}\n = C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y.\n\\end{equation}\n\\end{lem}\n\n\\begin{proof}\nBy Lemma~\\ref{L_4Y12_GetAuto},\nthere exists a function $x \\mapsto \\alpha_{x}$\nfrom $X$ to ${\\mathrm{Aut}} (D)$\nsuch that $x \\mapsto \\alpha_{x} (d)$ is continuous{} for all $d \\in D$\nand which satisfies $\\alpha (a) (x) = \\alpha_x (a (h^{-1} (x))$\nfor all $a \\in C_0 (X, D)$ and $x \\in X$.\n\nMost of the proof of Proposition~7.5 of~\\cite{PhLg}\ngoes through with only the obvious changes.\nIn analogy with that proof,\ndefine\n\\[\nB = \\big\\{ a \\in C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr) \\colon\n {\\mbox{$E_{\\alpha} (a u^{-n}) \\in C_0 (X \\setminus Y_n, \\, D)$\n for all $n \\in {\\mathbb{Z}}$}} \\big\\}\n\\]\nand\n\\[\nB_0 = B \\cap C (X) [{\\mathbb{Z}}].\n\\]\nThen $B_0$ is dense in~$B$ by the same reasoning as in~\\cite{PhLg}\n(using Ces\\`{a}ro means and Theorem VIII.2.2 of~\\cite{Dv}).\n\nThe proof of Proposition~7.5 of~\\cite{PhLg}\nshows that\nwhen $0 \\leq m \\leq n$ and also when $0 \\geq m \\geq n$,\nwe have $Y_m \\subset Y_n$,\nthat for $n \\in {\\mathbb{Z}}$,\nwe have\n\\begin{equation}\\label{Eq_4Y14_hnyn}\nh^{-n} (Y_n) = Y_{-n},\n\\end{equation}\nand that for $m, n \\in {\\mathbb{Z}}$,\nwe have $Y_{m + n} \\subset Y_m \\cup h^m (Y_n)$.\nIt then follows,\nas in~\\cite{PhLg},\nthat $B_0$ is a *-algebra,\nand that\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y\n \\subset {\\overline{B_0}} = B$.\n\nWe next claim that for all $n \\in {\\mathbb{Z}}$\nand $f \\in C_0 (X \\setminus Y_n, \\, D)$,\nwe have $f u^n \\in C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$.\nThe changes to the proof of Proposition~7.5 of~\\cite{PhLg}\nat this point are more substantial.\n\nFor $n = 0$ the claim is trivial.\nLet $n > 0$\nand let $f \\in C_0 (X \\setminus Y_n, \\, D)$.\nDefine $b = (f^* f)^{1 \/ (2 n)}$.\nLet $s \\in C_0 (X \\setminus Y_n, \\, D)''$\nbe the partial isometry in the polar decomposition of~$f$,\nso that $f = s (f^* f)^{1 \/ 2} = s b^n$.\nIt follows from Proposition~1.3 of~\\cite{Cu0}\nthat the element $a_0 = s b$ is in $C_0 (X \\setminus Y_n, \\, D)$.\nMoreover,\n$a_0 (f^* f)^{\\frac{1}{2} - \\frac{1}{2 n}} = f$.\nDefine $a_1 \\in C (X, D)$\nby $a_1 (x) = \\alpha_{h (x)}^{-1} \\big( b (h (x)) \\big)$\nfor $x \\in X$,\nand for $k = 1, 2, \\ldots, n - 2$ inductively define\n$a_{k + 1} \\in C (X, D)$\nby $a_{k + 1} (x) = \\alpha_{h (x)}^{-1} \\big( a_k (h (x)) \\big)$\nfor $x \\in X$.\nThe definition of $Y_n$\nimplies that\n$a_1, a_2, \\ldots, a_{n - 1} \\in C_0 (X \\setminus Y, \\, D)$,\nand we already have\n$a_0 \\in C_0 (X \\setminus Y_n, \\, D) \\S C_0 (X \\setminus Y, \\, D)$.\nTherefore the element\n\\[\na = (a_0 u) (a_1 u) \\cdots (a_{n - 1} u)\n\\]\nis in $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$.\nFor $k = 1, 2, \\ldots, n - 2$ and $x \\in X$, we have\n\\[\n\\alpha (a_{k + 1} ) (x)\n = \\alpha_x (a_{k + 1} (h^{-1} (x))\n = \\alpha_x \\big( \\alpha_{x}^{-1} ( a_k (x) ) \\big)\n = a_k (x),\n\\]\nso $\\alpha (a_{k + 1} ) = a_k$.\nSimilarly, $\\alpha (a_1) = b$.\nNow\n\\begin{align*}\na\n& = a_0 (u a_1 u^{-1}) (u^2 a_2 u^{-2})\n \\cdots \\big( u^{n - 1} a_{n - 1} u^{- (n - 1)} \\big) u^n\n \\\\\n& = a_0 \\alpha (a_1) \\alpha^2 (a_2)\n \\cdots \\alpha^{n - 1} (a_{n - 1}) u^n\n = (s b) b^{n - 1} u^n\n = f u^n.\n\\end{align*}\nSo $f u^n \\in C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$.\nFinally,\nsuppose $n < 0$,\nand let $f \\in C_0 (X \\setminus Y_n, \\, D)$.\nIt follows from~(\\ref{Eq_4Y14_hnyn})\nthat $f \\circ h^n \\in C_0 (X \\setminus Y_{- n}, \\, D)$,\nwhence also $(f \\circ h^n)^* \\in C_0 (X \\setminus Y_{- n}, \\, D)$.\nSince $- n > 0$,\nwe therefore get\n\\[\nf u^n = \\big( u^{-n} f^* \\big)^*\n = \\big( (f \\circ h^n)^* u^{- n} \\big)^*\n \\in C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y.\n\\]\nThe claim is proved.\n\nIt now follows that\n$B_0 \\subset C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$.\nCombining this result with ${\\overline{B_0}} = B$\nand $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y \\subset B$,\nwe get $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y = B$.\n\\end{proof}\n\nThe following lemma is the analog of Lemma~7.8 of~\\cite{PhLg}.\n\n\\begin{lem}\\label{L_4X21_78}\nLet $G$ be a discrete group,\nlet $X$ be a compact space,\nand suppose $G$ acts on $X$ in such a way that\nfor every finite set $S \\subset G \\setminus \\{ 1 \\}$,\nthe set\n\\[\n\\big\\{ x \\in X \\colon {\\mbox{$g x \\neq x$ for all $g \\in S$}} \\big\\}\n\\]\nis dense in~$X$.\nLet $D$ be a unital C*-algebra,\nand let $\\alpha \\colon G \\to {\\mathrm{Aut}} (C (X, D))$\nbe an action of $G$ on $C (X, D)$\nwhich lies over the given action of $G$ on $X$\n(in the sense of Definition~\\ref{D_4Y12_OverX}).\nFollowing Notation~\\ref{N_4Y12_CPN},\nlet\n$a \\in C (X, D) [G]\n \\subset C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nand let $\\varepsilon > 0$.\nThen there exists $f \\in C (X) \\subset C (X, D)$\nsuch that\n\\[\n0 \\leq f \\leq 1,\n\\qquad\nf a^* a f \\in C (X, D),\n\\qquad {\\mbox{and}} \\qquad\n\\| f a^* a f \\| \\geq \\| E_{\\alpha} (a^* a) \\| - \\varepsilon.\n\\]\n\\end{lem}\n\n\\begin{proof}\nWe follow the proof of Lemma~7.8 of~\\cite{PhLg}.\nSet $b = a^* a$.\nThere are a finite set $T \\subset G$\nwith $1 \\in T$,\nand elements $b_g \\in C (X, D)$ for $g \\in T$,\nsuch that $b = \\sum_{g \\in T} b_g u_g$.\nIf $\\| b_1 \\| \\leq \\varepsilon$ take $f = 0$.\nOtherwise, define\n\\[\nU = \\big\\{ x \\in X \\colon\n \\| b_1 (x) \\| > \\| E (a^* a) \\| - \\varepsilon \\big\\},\n\\]\nwhich is a nonempty open subset of~$X$.\nChoose $V$, $W$, $f$, and $x_0$\nas in the proof of Lemma~7.8 of~\\cite{PhLg}.\nThen,\nas there, $f b f = f b_1 f$.\nMoreover,\n\\[\n\\| f b_1 f \\|\n \\geq f (x_0) \\| b_1 (x_0) \\| f (x_0)\n = \\| b_1 (x_0) \\|\n > \\| E_{\\alpha} (a^* a) \\| - \\varepsilon.\n\\]\nThis completes the proof.\n\\end{proof}\n\nThe following lemma is the analog of Lemma~7.9 of~\\cite{PhLg}.\n\n\\begin{lem}\\label{L_4X21_79}\nLet $G$, $X$, the action of $G$ on $X$,\n$D$, and $\\alpha \\colon G \\to {\\mathrm{Aut}} (C (X, D))$\nbe as in Lemma~\\ref{L_4X21_78}.\nLet $B \\subset C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nbe a unital subalgebra\nsuch that,\nfollowing Notation~\\ref{N_4Y12_CPN},\n\\begin{enumerate}\n\\item\\label{L_4X21_79:CXI}\n$C (X, D) \\subset B$.\n\\item\\label{L_4X21_79:FSD}\n$B \\cap C (X, D)[G]$ is dense in~$B$.\n\\end{enumerate}\nLet $a \\in B_{+} \\setminus \\{ 0 \\}$.\nThen there exists $b \\in C (X, D)_{+} \\setminus \\{ 0 \\}$\nsuch that $b \\precsim_{B} a$.\n\\end{lem}\n\n\\begin{proof}\nWe follow the proof of Lemma~7.9 of~\\cite{PhLg},\nusing our Lemma~\\ref{L_4X21_78}\nin place of Lemma~7.8 of~\\cite{PhLg},\nexcept that $c \\in B \\cap C (X, D)[G]$.\nThe element $(f c^* c f - 2 \\varepsilon)_{+}$\nwe obtain now satisfies\n$(f c^* c f - 2 \\varepsilon)_{+} \\in C (X, D)_{+} \\setminus \\{ 0 \\}$\nand $(f c^* c f - 2 \\varepsilon)_{+} \\precsim_{B} a$.\n\\end{proof}\n\nIn Lemma~\\ref{L_4X21_79},\nwe really want to have $b \\in C (X)_{+} \\setminus \\{ 0 \\}$.\nWhen $G = {\\mathbb{Z}}$ and under the pseudoperiodicity hypothesis\nof \\Def{D_4Y16_PsPrGen}, this is possible.\n\n\\begin{lem}\\label{L_4Y11_Comp}\nLet $X$ be a compact metric space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a simple unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nAssume that the action generated by $\\alpha$\nis pseudoperiodically generated.\nLet $Y \\subset X$ be a compact set\nsuch that $h^n (Y) \\cap Y = \\varnothing$\nfor all $n \\in {\\mathbb{Z}} \\setminus \\{ 0 \\}$.\nSet $B = C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$.\nThen for every $a \\in C (X, D)_{+} \\setminus \\{ 0 \\}$\nthere exists $f \\in C (X)_{+} \\setminus \\{ 0 \\} \\subset B$\nsuch that $f \\precsim_B a$.\n\\end{lem}\n\n\\begin{proof}\nUse Kirchberg's Slice Lemma (Lemma 4.1.9 of~\\cite{Rrd})\nto find $g_0 \\in C (X)_{+} \\setminus \\{ 0 \\}$\nand $d_0 \\in D_{+} \\setminus \\{ 0 \\}$\nsuch that\n\\begin{equation}\\label{Eq_5417_Zero}\ng_0 \\otimes d_0 \\precsim_{C (X, D)} a.\n\\end{equation}\nAs in Definition~\\ref{D_4Y12_OverX},\nlet $(m, x) \\mapsto \\alpha_{m, x}$\nbe the function ${\\mathbb{Z}} \\times X \\to {\\mathrm{Aut}} (D)$\nsuch that $\\alpha^m (a) (x) = \\alpha_{m, x} (a (h^{-m} (x))$\nfor $m \\in {\\mathbb{Z}}$, $a \\in C_0 (X, D)$, and $x \\in X$.\nSince the action generated by $\\alpha$\nis pseudoperiodically generated,\nthere exists $d_1\\in D_{+} \\setminus \\{ 0 \\}$\nsuch that\n\\begin{equation}\\label{Eq_5418_PsPer}\n\\alpha_{m, x} (d_1) \\precsim d_0\n\\end{equation}\nfor all $x \\in X$ and $m \\in {\\mathbb{Z}}$.\nWithout loss of generality{} $\\| d_1 \\| = 1$.\nSet $d = \\big( d_1 - \\tfrac{1}{2} \\big)_{+}$.\nUse Corollary 1.14 of~\\cite{PhLg}\nto find $k \\in {\\mathbb{Z}}_{> 0}$ and $w_1, w_2, \\ldots, w_k \\in D$\nsuch that\n\\begin{equation}\\label{Eq_5418_Sum1}\n\\sum_{j = 1}^k w_j d w_j^* = 1.\n\\end{equation}\n\nThe set $X \\setminus Y$ is dense in~$X$,\nso there is $x_0 \\in X \\setminus Y$ such that $g_0 (x_0) \\neq 0$.\nChoose $g_1 \\in C (X)_{+}$ such that $g_1 (x_0) = 1$\nand $g_1 |_Y = 0$.\nThen\n\\begin{equation}\\label{Eq_5417_One}\ng_1 g_0 \\otimes d_0\n \\precsim_{C (X, D)} g_0 \\otimes d_0.\n\\end{equation}\nSet\n$U_0 = \\big\\{ x \\in X \\colon (g_1 g_0) (x) \\neq 0 \\big\\}$.\nChoose a nonempty open set $U \\subset X$\nsuch that ${\\overline{U}} \\subset U_0$.\nSet\n\\begin{equation}\\label{Eq_9Z01_NewLabel}\n\\rho = \\inf_{x \\in {\\overline{U}}} (g_0 g_1) (x).\n\\end{equation}\nThen $\\rho > 0$.\nThe set\n\\[\nX \\setminus \\bigcup_{n \\in {\\mathbb{Z}}} h^n (Y)\n = \\bigcap_{n \\in {\\mathbb{Z}}} h^n (X \\setminus Y)\n\\]\nis dense by the Baire Category Theorem.\nSo we can choose\n$y_0 \\in U \\cap \\left( X \\setminus\n\\bigcup_{n \\in {\\mathbb{Z}}} h^n (Y) \\right)$.\nThe forward orbits $\\{ h^n (y_0) \\colon n \\geq N \\}$\nof $y_0$ are all dense,\nso there exist $n_1, n_2, \\ldots, n_k \\in {\\mathbb{Z}}_{\\geq 0}$\nwith $0 = n_1 < n_2 < \\cdots< n_k$\nand $h^{n_j} (y_0) \\in U$ for $j = 1, 2, \\ldots, k$.\n\nChoose an open set $V$ containing $y_0$ which is so small that\nthe following hold:\n\\begin{enumerate}\n\\item\\label{L_4Y11_Comp_Cont}\n${\\overline{V}} \\subset X \\setminus\n\\bigcup_{m = 0}^{n_k} h^{- m} (Y)$.\n\\item\\label{L_4Y11_Comp_InU}\n$h^{n_j} ( {\\overline{V}} ) \\subset U$\nfor $j = 1, 2, \\ldots, k$.\n\\item\\label{L_4Y11_Comp_Disj}\nThe sets\n${\\overline{V}}, \\, h ( {\\overline{V}} ), \\, \\ldots, \\,\n h^{n_k} ( {\\overline{V}} )$\nare disjoint.\n\\item\\label{L_4Y11_Comp_Small}\nFor $y \\in V$ and $m = 0, 1, \\ldots, n_k$,\nwe have\n\\[\n\\bigl\\| \\alpha_{m, h^m (y)} (d_1) - \\alpha_{m, h^m (y_0)} (d_1) \\bigr\\|\n< \\frac{1}{4}.\n\\]\n\\setcounter{TmpEnumi}{\\value{enumi}}\n\\end{enumerate}\nChoose $f_0, f \\in C (X)_{+}$\nsuch that:\n\\begin{enumerate}\n\\setcounter{enumi}{\\value{TmpEnumi}}\n\\item\\label{L_4Y11_Comp_1}\n$\\| f_0 \\|, \\, \\| f \\| \\leq 1$.\n\\item\\label{L_4Y11_Comp_Supp}\n${\\mathrm{supp}} (f_0) \\subset V$.\n\\item\\label{L_4Y11_Comp_b0b}\n$f_0 f = f$.\n\\item\\label{L_4Y11_Comp_Aty0}\n$f (y_0) = 1$.\n\\end{enumerate}\nFor $m = 0, 1, \\ldots, n_k - 1$,\ndefine $c_m \\in C (X, D)$\nby $c_m (x) = f ( h^{- m} (x)) \\alpha_{m, x} (d)$\nfor $x \\in X$.\nThus\n\\begin{equation}\\label{Eq_5417_cm}\nc_m = \\alpha^m (f \\otimes d).\n\\end{equation}\nAlso set\n\\begin{equation}\\label{Eq_5418_DefBt}\n\\beta_m = \\alpha_{m, h^m (y_0)}.\n\\end{equation}\n\nWe claim that for $m = 0, 1, \\ldots, n_k$\nwe have\n\\[\nc_m \\precsim_{C (X, D)} (f \\circ h^{-m}) \\otimes \\beta_m (d_1).\n\\]\n\nTo prove the claim, let $\\varepsilon > 0$.\nSet $L = h^{m} ( {\\mathrm{supp}} (f_0) )$.\nDefine $b \\in C (L, D)$\nby $b (x) = \\alpha_{m, x} (d_1)$ for $x \\in L$.\nSince $L \\subset h^m (V)$,\ncondition~(\\ref{L_4Y11_Comp_Small})\nimplies that\nfor $x \\in L$\nwe have $\\| \\alpha_{m, x} (d_1) - \\beta_m (d_1) \\| < \\frac{1}{2}$.\nIn $C (L, D)$ we therefore get\n$\\| b - 1 \\otimes \\beta_m (d_1) \\|\n \\leq \\frac{1}{4} < \\frac{1}{2}$.\nSo\n\\[\n\\bigl( b - \\tfrac{1}{2} \\bigr)_{+}\n \\precsim_{C (L, D)} 1 \\otimes \\beta_m (d_1).\n\\]\nTherefore there exists $v_0 \\in C (L, D)$\nsuch that\n\\[\n\\big\\| v_0^* ( 1 \\otimes \\beta_m (d_1) ) v_0\n - \\bigl( b - \\tfrac{1}{2} \\bigr)_{+} \\big\\|\n < \\frac{\\varepsilon}{2}.\n\\]\nDefine $v, a \\in C (X, D)$ by\n\\[\nv (x) = \\begin{cases}\n f_0 (h^{- m} (x)) v_0 (x) & x \\in L\n \\\\\n 0 & x \\not\\in L\n\\end{cases}\n\\qquad {\\mbox{and}} \\qquad\na = v^* \\bigl[ (f \\circ h^{-m}) \\otimes \\beta_m (d_1) \\bigr] v.\n\\]\nWe show that $\\| a - c_m \\| < \\varepsilon$.\n\nFor $x \\in X \\setminus L$,\nwe have $a (x) = c_m (x) = 0$.\nFor $x \\in L$,\nusing $f_0 f = f$ at the second step\nand using $d = \\bigl( d_1 - \\tfrac{1}{2} \\bigr)_{+}$\nand the definition of~$b$ at the third step,\nwe get\n\\begin{align*}\n& \\| a (x) - c_m (x) \\|\n\\\\\n& \\hspace*{3em} {\\mbox{}}\n = \\big\\| f_0 (h^{- m} (x))^2 f (h^{- m} (x)) v_0 (x)^*\n \\beta_m (d_1) v_0 (x)\n - f (h^{- m} (x)) \\alpha_{m, x} (d) \\big\\|\n\\\\\n& \\hspace*{3em} {\\mbox{}}\n = \\big\\| f (h^{- m} (x))\n [ v_0 (x)^* \\beta_m (d_1) v_0 (x) - \\alpha_{m, x} (d) ] \\big\\|\n\\\\\n& \\hspace*{3em} {\\mbox{}}\n = f (h^{- m} (x))\n \\big\\| v_0 (x)^* \\beta_m (d_1) v_0 (x)\n - \\bigl( b (x) - \\tfrac{1}{2} \\bigr)_{+} \\big\\|\n\\\\\n& \\hspace*{3em} {\\mbox{}}\n \\leq f (h^{- m} (x))\n \\big\\| v_0^* ( 1 \\otimes \\beta_m (d_1) ) v_0\n - \\bigl( b - \\tfrac{1}{2} \\bigr)_{+} \\big\\|\n < \\frac{\\varepsilon}{2}.\n\\end{align*}\nTaking the supremum over $x \\in X$ gives\n$\\| a - c_m \\| \\leq \\frac{\\varepsilon}{2} < \\varepsilon$.\nThus\n\\[\n\\big\\| v^* [ (f \\circ h^{-m}) \\otimes \\beta_m (d_1) ] v - c_m \\big\\|\n < \\varepsilon.\n\\]\nSince $\\varepsilon > 0$ is arbitrary,\nthe claim follows.\n\nFor $m = 0, 1, \\ldots, n_k$,\nthe functions $f \\circ h^{-m}$ are orthogonal\nsince the sets $h^m ( {\\overline{V}} )$ are disjoint.\nThe claim therefore implies that\n\\begin{equation}\\label{Eq_5417_Two}\n\\sum_{j = 1}^k c_{n_j}\n \\precsim_{C (X, D)}\n \\sum_{j = 1}^k (f \\circ h^{-n_j}) \\otimes \\beta_{n_j} (d_1).\n\\end{equation}\n\nWe now claim that\n\\begin{equation}\\label{Eq_5417_Three}\nf \\otimes 1 \\precsim_{B} \\sum_{j = 1}^k c_{n_j}.\n\\end{equation}\nTo prove this claim,\nfor $j = 1, 2, \\ldots, k$\ndefine\n\\[\nv_j = [(f \\circ h^{-n_j}) \\otimes 1] u^{n_j} (1 \\otimes w_j)^*\n \\in C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr).\n\\]\nCombining (\\ref{L_4Y11_Comp_Supp}) and~(\\ref{L_4Y11_Comp_Cont}),\nwe see that $f_0$ vanishes in particular on the sets\n\\[\nh^{-1} ( Y ), \\, h^{-2} ( Y ), \\, \\ldots, \\, h^{- n_j} ( Y ),\n\\]\nwhence $f_0 \\circ h^{- n_j}$ vanishes on the sets\n\\[\nY, \\, h ( Y ), \\, \\ldots, \\, h^{n_j - 1} ( Y ).\n\\]\nSo $v_j \\in B$ by\n\\Lem{L_4X21_AYStruct}.\nUsing (\\ref{Eq_5417_cm}) at the first step\nand $f_0 f = f$ at the last step,\nwe calculate:\n\\begin{align*}\nv_j^* c_{n_j} v_j\n& = (1 \\otimes w_j) u^{- n_j} [(f \\circ h^{-n_j}) \\otimes 1]\n \\alpha^{n_j} (f \\otimes d) [(f \\circ h^{-n_j}) \\otimes 1]\n u^{n_j} (1 \\otimes w_j)^*\n\\\\\n& = (1 \\otimes w_j) (f_0 \\otimes 1) (f \\otimes d)\n (f_0 \\otimes 1) (1 \\otimes w_j^*)\n = f \\otimes w_j d w_j^*.\n\\end{align*}\nWe apply (\\ref{Eq_5418_Sum1}) at the first step\nand use orthogonality\nof $c_0, c_1, \\ldots, c_{n_k}$\nand Lemma 1.4(12) of~\\cite{PhLg} at the second step,\nto get\n\\[\nf \\otimes 1\n = \\sum_{j = 1}^k f \\otimes w_j d w_j^*\n \\precsim_B \\sum_{j = 1}^k c_{n_j}.\n\\]\nThis proves the claim~(\\ref{Eq_5417_Three}).\n\nWe next claim that\n\\begin{equation}\\label{Eq_5417_Four}\n\\sum_{j = 1}^k (f \\circ h^{-n_j}) \\otimes \\beta_{n_j} (d_1)\n \\precsim_{C (X, D)} g_0 g_1 \\otimes d.\n\\end{equation}\nTo prove this,\ncombine (\\ref{Eq_5418_PsPer}), (\\ref{Eq_5418_DefBt}),\nand Lemma 1.11 of~\\cite{PhLg}\nto get\n\\[\n(f \\circ h^{-n_j}) \\otimes \\beta_{n_j} (d_1)\n \\precsim_{C (X, D)} (f \\circ h^{-n_j}) \\otimes d_0\n\\]\nfor $j = 1, 2, \\ldots, k$.\nSince the functions\n$f \\circ h^{-n_j}$ are orthogonal,\nit follows that\n\\[\n\\sum_{j = 1}^k (f \\circ h^{-n_j}) \\otimes \\beta_{n_j} (d_1)\n \\precsim_{C (X, D)}\n \\left( \\sum_{j = 1}^k (f \\circ h^{-n_j}) \\right) \\otimes d_0.\n\\]\nUsing orthogonality\nof the functions\n$f \\circ h^{-n_j}$ again,\ntogether with\n\\[\n{\\mathrm{supp}} (f \\circ h^{-n_j}) \\subset h^{n_j} ( {\\overline{V}} )\n \\subset U\n\\qquad {\\mbox{and}} \\qquad\n0 \\leq f \\circ h^{-n_j} \\leq 1\n\\]\nfor $j = 1, 2, \\ldots, k$,\nwe see that\n$0 \\leq \\sum_{j = 1}^k f \\circ h^{-n_j} \\leq 1$\nand,\nby~(\\ref{L_4Y11_Comp_InU}), that this sum is supported in~$U$.\nSince $g_0 g_1 \\geq \\rho$ on ${\\overline{U}}$\nby~(\\ref{Eq_9Z01_NewLabel}),\nwe get the first step in the following computation;\nthe second step is clear:\n\\[\n\\sum_{j = 1}^k (f \\circ h^{-n_j}) \\otimes d_0\n \\leq \\frac{1}{\\rho} ( g_0 g_1 \\otimes d_0 )\n \\sim_{C (X, D)} g_0 g_1 \\otimes d_0.\n\\]\nThe claim is proved.\n\nNow combine\n(\\ref{Eq_5417_Zero}),\n(\\ref{Eq_5417_One}),\n(\\ref{Eq_5417_Four}),\n(\\ref{Eq_5417_Two}),\nand\n(\\ref{Eq_5417_Three})\nto get $f \\otimes 1 \\precsim_{B} a$.\n\\end{proof}\n\n\\begin{cor}\\label{C_5418_CompCX}\nLet $X$ be a compact metric space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a simple unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nAssume that the action generated by $\\alpha$\nis pseudoperiodically generated.\nLet $Y \\subset X$ be a compact set\n(possibly empty)\nsuch that $h^n (Y) \\cap Y = \\varnothing$\nfor all $n \\in {\\mathbb{Z}} \\setminus \\{ 0 \\}$.\nSet $B = C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$.\nLet $a \\in B_{+} \\setminus \\{ 0 \\}$.\nThen there exists $f \\in C (X)_{+} \\setminus \\{ 0 \\} \\subset B$\nsuch that $f \\precsim_B a$.\n\\end{cor}\n\n\\begin{proof}\n\\Lem{L_4X21_AYStruct}\nimplies that $B$ satisfies the hypotheses\nof \\Lem{L_4X21_79}\nwith $G = {\\mathbb{Z}}$.\nTherefore,\nby \\Lem{L_4X21_79},\nthere exists $b \\in C (X, D)_{+} \\setminus \\{ 0 \\}$\nsuch that $b \\precsim_{B} a$.\nNow \\Lem{L_4Y11_Comp}\nprovides $f \\in C (X)_{+} \\setminus \\{ 0 \\}$\nsuch that $f \\precsim_B a$.\n\\end{proof}\n\nWe can now prove the analog of Theorem~7.10 of~\\cite{PhLg}.\n\n\\begin{thm}\\label{T_5418_AYLg}\nLet $X$ be a compact metric space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a simple unital C*-algebra{} which has a tracial state,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nAssume that the action generated by $\\alpha$\nis pseudoperiodically generated.\nLet $Y \\subset X$ be a compact subset\nsuch that $h^n (Y) \\cap Y = \\varnothing$\nfor all $n \\in {\\mathbb{Z}} \\setminus \\{ 0 \\}$.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis a large subalgebra of $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nof crossed product type\nin the sense of Definition~4.9 of~\\cite{PhLg}.\n\\end{thm}\n\n\\begin{proof}\nWe verify the hypotheses of\nProposition~4.11 of~\\cite{PhLg}.\nWe follow Notation~\\ref{N_4Y12_CPN}.\nSet\n\\[\nA = C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr),\n\\qquad\nB = C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y,\n\\]\n\\[\nC = C (X, D),\n\\qquad {\\mbox{and}} \\qquad\nG = \\{ u \\}.\n\\]\n\nThe algebra $A$ is simple\nby Proposition~\\ref{P_4Y15_CPSimple}\nand finite by Proposition~\\ref{P_5418_CPStFin}.\nIn particular,\ncondition~(1)\nin Proposition 4.11 of~\\cite{PhLg} holds.\n\nWe next verify condition~(2)\nin Proposition 4.11 of~\\cite{PhLg}.\nAll parts are obvious except~(2d).\nSo let $a \\in A_{+} \\setminus \\{ 0 \\}$\nand $b \\in B_{+} \\setminus \\{ 0 \\}$.\nApply Corollary~\\ref{C_5418_CompCX} twice,\nthe first time with $Y = \\varnothing$ and $a$ as given\nand the second time with $Y$ as given\nand with $b$ in place of~$a$.\nWe get $a_0, b_0 \\in C (X)_{+} \\setminus \\{ 0 \\}$\nsuch that $a_0 \\precsim_{A} a$ and $b_0 \\precsim_{B} b$.\n\nWe can now argue as in the corresponding part of the\nproof of Theorem~7.10 of~\\cite{PhLg}\nto find $f \\in C (X)$\nsuch that\n\\[\nf \\precsim_{C (X)} a_0 \\precsim_{A} a\n\\qquad {\\mbox{and}} \\qquad\nf \\precsim_{B} b_0 \\precsim_{B} b,\n\\]\ncompleting the proof of condition~(2d).\nAlternatively,\n$a_0, b_0\n \\in C^* ({\\mathbb{Z}}, X, h)_Y\n \\subset C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$,\nand $C^* ({\\mathbb{Z}}, X, h)_Y$ is a large subalgebra of $C^* ( {\\mathbb{Z}}, X, h)$\nof crossed product type\nby Theorem~7.10 of~\\cite{PhLg},\nso the existence of $f$\nfollows from condition~(2d)\nin Proposition 4.11 of~\\cite{PhLg}.\n\nWe next verify condition~(3)\nin Proposition 4.11 of~\\cite{PhLg}.\nLet $m \\in {\\mathbb{Z}}_{> 0}$,\nlet $a_1, a_2, \\ldots, a_m \\in A$,\nlet $\\varepsilon > 0$,\nand let $b \\in B_{+} \\setminus \\{ 0 \\}$.\nWe follow the corresponding part of the\nproof of Theorem~7.10 of~\\cite{PhLg}.\nChoose $c_1, c_2, \\ldots, c_m \\in C (X, D) [{\\mathbb{Z}}]$\nsuch that $\\| c_j - a_j \\| < \\varepsilon$\nfor $j = 1, 2, \\ldots, m$.\n(This estimate is condition~(3b).)\nChoose $N \\in {\\mathbb{Z}}_{> 0}$\nsuch that for $j = 1, 2, \\ldots, m$\nthere are $c_{j, l} \\in C (X, D)$\nfor $l = -N, \\, - N + 1, \\, \\ldots, \\, N - 1, \\, N$\nwith\n\\[\nc_j = \\sum_{l = -N}^N c_{j, l} u^l.\n\\]\n\nApply Corollary~\\ref{C_5418_CompCX} to~$B$,\nto find $f \\in C (X)_{+} \\setminus \\{ 0 \\}$\nsuch that $f \\precsim_{B} b$.\nSet $U = \\{ x \\in X \\colon f (x) \\neq 0 \\}$,\nand choose nonempty disjoint open sets\n$U_l \\subset U$ for $l = -N, \\, - N + 1, \\, \\ldots, \\, N - 1, \\, N$.\nFor each such~$l$,\nuse Lemma~7.7 of~\\cite{PhLg}\nto choose $f_l, r_l \\in C (X)_{+}$\nsuch that $r_l (x) = 1$ for all $x \\in h^l (Y)$,\nsuch that $0 \\leq r_l \\leq 1$,\nsuch that ${\\mathrm{supp}} (f_l) \\subset U_l$,\nand such that $r_l \\precsim_{C^* ({\\mathbb{Z}}, X, h)_Y} f_l$.\nThen also $r_l \\precsim_{B} f_l$.\n\nChoose an open set $W$ containing~$Y$\nsuch that\n\\[\nh^{-N} (W), \\, h^{- N + 1} (W),\n \\, \\ldots, \\, h^{N - 1} (W), \\, h^{N} (W)\n\\]\nare disjoint,\nand choose $r \\in C (X)$ such that $0 \\leq r \\leq 1$,\n$r (x) = 1$ for all $x \\in Y$,\nand ${\\mathrm{supp}} (r) \\subset W$.\nSet\n\\[\ng_0 = r \\cdot \\prod_{l = - N}^{N} r_l \\circ h^{l}.\n\\]\nSet $g_l = g_0 \\circ h^{- l}$\nfor $l = -N, \\, - N + 1, \\, \\ldots, \\, N - 1, \\, N$.\nThen $0 \\leq g_l \\leq r_l \\leq 1$.\nSet $g = \\sum_{l = - N}^{N} g_l$.\nThe supports of the functions $g_l$ are disjoint,\nso $0 \\leq g \\leq 1$.\nThis is condition~(3a)\nin Proposition 4.11 of~\\cite{PhLg}.\nThe proof of Theorem~7.10 of~\\cite{PhLg}\nshows that\n$g \\precsim_{C^* ({\\mathbb{Z}}, X, h)_Y} f \\precsim_B b$.\nSince $C^* ({\\mathbb{Z}}, X, h)_Y \\subset B$,\nit follows that $g \\precsim_B b$.\nThis is condition~(3d)\nin Proposition 4.11 of~\\cite{PhLg}.\n\nIt remains to verify condition~(3c)\nin Proposition 4.11 of~\\cite{PhLg}.\nThis is done the same way as\nthe corresponding part of the proof of Theorem~7.10 of~\\cite{PhLg},\nexcept using \\Lem{L_4X21_AYStruct}\nin place of Proposition~7.5 of~\\cite{PhLg}.\n\\end{proof}\n\nThe following corollary is the analog of Corollary~7.11\nof~\\cite{PhLg}.\n\n\\begin{cor}\\label{C_5418_AYStabLg}\nLet $X$ be a compact metric space, let $h \\colon X \\to X$ be a minimal\nhomeomorphism, let $D$ be a simple unital C*-algebra{}\nwhich has a tracial state,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nAssume that the action generated by $\\alpha$\nis pseudoperiodically generated.\nLet $Y \\subset X$ be a compact subset\nsuch that $h^n (Y) \\cap Y = \\varnothing$\nfor all $n \\in {\\mathbb{Z}} \\setminus \\{ 0 \\}$.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis a stably large subalgebra\nof $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nin the sense of Definition~5.1 of~\\cite{PhLg}.\n\\end{cor}\n\n\\begin{proof}\nSince $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is stably finite\n(by Proposition~\\ref{P_5418_CPStFin}),\nwe can combine Theorem~\\ref{T_5418_AYLg},\nProposition~4.10 of~\\cite{PhLg},\nand Corollary~5.8 of~\\cite{PhLg}.\n\\end{proof}\n\n\\begin{cor}\\label{AyCentLarge}\nLet $X$ be a compact metric space, let $h \\colon X \\to X$ be a minimal\nhomeomorphism, let $D$ be a simple unital C*-algebra{}\nwhich has a tracial state,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nAssume that the action generated by $\\alpha$\nis pseudoperiodically generated.\nLet $Y \\subset X$ be a compact subset\nsuch that $h^n (Y) \\cap Y = \\varnothing$\nfor all $n \\in {\\mathbb{Z}} \\setminus \\{ 0 \\}$.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis a centrally large subalgebra\nof $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nin the sense of Definition~3.2 of~\\cite{ArPh}.\n\\end{cor}\n\n\\begin{proof}\nSince $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is stably finite\n(by Proposition~\\ref{P_5418_CPStFin}),\nwe can combine Theorem~\\ref{T_5418_AYLg}\nand Theorem~4.6 of~\\cite{ArPh}.\n\\end{proof}\n\nWe conclude by giving some conditions on $D$ and~$\\alpha$ which\nguarantee the hypotheses of Corollary~\\ref{AyCentLarge}.\nThese are more natural to consider than the awkward\npseudoperiodicity hypothesis.\n\n\\begin{cor}\\label{AyCentLargeConds}\nLet $X$ be a compact metric space,\nlet $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $Y \\subset X$ be a compact subset\nsuch that $h^{n} (Y) \\cap Y = \\varnothing$\nfor all $n \\in {\\mathbb{Z}} \\setminus \\set{0}$,\nlet $D$ be a simple unital C*-algebra{}\nwhich has a tracial state,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nLet $x \\mapsto \\alpha_{x}$\nbe the corresponding map from $X$ to ${\\mathrm{Aut}} (D)$,\nas in Lemma~\\ref{L_4Y12_GetAuto}.\nAssume one of the following conditions holds.\n\\begin{enumerate}\n\\item\\label{9212_AyCentLargeConds_AppInn}\nAll elements of $\\set{ \\alpha_{x} \\colon x \\in X} \\subset {\\mathrm{Aut}} (D)$\nare approximately inner.\n\\item\\label{9212_AyCentLargeConds_StrCmp}\n$D$ has strict comparison of positive elements.\n\\item\\label{9212_AyCentLargeConds_SP}\n$D$ has property (SP) and the order on projections over~$D$\nis determined by quasitraces.\n\\item\\label{9212_AyCentLargeConds_CptGen}\nThe set\n$\\set{ \\alpha_{x} \\colon x \\in X}$\nis contained in a subgroup of ${\\mathrm{Aut}} (D)$\nwhich is compact in the topology\nof pointwise convergence in the norm on~$D$.\n\\end{enumerate}\nThen\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis a centrally large subalgebra of\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$.\n\\end{cor}\n\n\\begin{proof}\nWe claim that, under any of the conditions\n(\\ref{9212_AyCentLargeConds_AppInn}),\n(\\ref{9212_AyCentLargeConds_StrCmp}),\n(\\ref{9212_AyCentLargeConds_SP}),\nor (\\ref{9212_AyCentLargeConds_CptGen}),\nthe set $\\set{ \\alpha_{x} \\colon x \\in X}$\nis contained in a pseudoperiodic subgroup of ${\\mathrm{Aut}} (D)$.\nFor~(\\ref{9212_AyCentLargeConds_AppInn}),\nthis follows from Lemma~\\ref{L_4Y12_ApproxInn};\nfor~(\\ref{9212_AyCentLargeConds_StrCmp}),\nfrom Lemma~\\ref{L_4Y12_TrPrsv};\nfor~(\\ref{9212_AyCentLargeConds_SP}),\nfrom Lemma~\\ref{L_5417_SP};\nand for~(\\ref{9212_AyCentLargeConds_CptGen}),\nfrom Lemma~\\ref{Lem1_190201D}.\nNow apply Lemma~\\ref{L_4Y16_WhenPPG}\nto see that\nthe hypotheses of Corollary~\\ref{AyCentLarge} are satisfied.\n\\end{proof}\n\nThe case in which $D$ is purely infinite and simple is\nmuch easier.\nWe give a direct proof, not depending on large subalgebras,\nthat $C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nis purely infinite simple for any discrete group~$G$.\nIt is still true,\nand will be proved below (Proposition \\ref{T_9921_PICase}),\nthat, under the other hypotheses of Theorem~\\ref{T_5418_AYLg}\nthe subalgebra $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis a large subalgebra of $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nof crossed product type.\nThis fact seems potentially useful,\nbut does not help with the analysis of any of the examples\nin Section~\\ref{Sec_InfRed}.\n\n\\begin{thm}\\label{T_0101_PICase}\nLet $G$ be a discrete group,\nlet $X$ be a compact space,\nand suppose $G$ acts on $X$ in such a way that\nthe action is minimal and\nfor every finite set $S \\subset G \\setminus \\{ 1 \\}$,\nthe set\n\\[\n\\big\\{ x \\in X \\colon {\\mbox{$g x \\neq x$ for all $g \\in S$}} \\big\\}\n\\]\nis dense in~$X$.\nLet $D$ be a purely infinite simple unital C*-algebra,\nand let $\\alpha \\colon G \\to {\\mathrm{Aut}} (C (X, D))$\nbe an action of $G$ on $C (X, D)$\nwhich lies over the given action of $G$ on $X$\n(in the sense of Definition~\\ref{D_4Y12_OverX}).\nThen $C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nis purely infinite simple.\n\\end{thm}\n\nWe saw in Proposition~\\ref{P_4Y15_CPSimple}\nthat $C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$\nis simple,\nbut we won't use that in this proof.\n\n\\begin{proof}[Proof of Theorem~\\ref{T_0101_PICase}]\nFor convenience of notation,\nset $A = C^*_{\\mathrm{r}} \\big( G, \\, C (X, D), \\, \\alpha \\big)$.\nWe also freely identify $C (X, D)$ with $C (X) \\otimes D$.\nWe prove the result by showing that $1_A \\precsim a$\nfor all $a \\in A_{+} \\setminus \\{ 0 \\}$.\nBy Lemma~\\ref{L_4X21_79}\n(taking $B$ there to be~$A$),\nwe can assume that $a \\in C (X, D)_{+} \\setminus \\{ 0 \\}$.\n\nUse Kirchberg's Slice Lemma (Lemma 4.1.9 of~\\cite{Rrd})\nto find $f \\in C (X)_{+} \\setminus \\{ 0 \\}$\nand $b \\in D_{+} \\setminus \\{ 0 \\}$\nsuch that\n\\begin{equation}\\label{Eq_0201_Slice}\nf \\otimes b \\precsim_{C (X, D)} a.\n\\end{equation}\nWithout loss of generality{} $\\| f \\| = 1$.\nSince $D$ is purely infinite and simple,\nwe have $1_D \\precsim_{D} b$, and so it follows that\n\\begin{equation}\\label{Eq_0201_InD}\nf \\otimes 1_D \\precsim_{D} f \\otimes b.\n\\end{equation}\nSet $U = \\bigl\\{ x \\in X \\colon f (x) > \\frac{1}{2} \\bigr\\}$.\nBy minimality of the action,\nthe sets $g U$ for $g \\in G$ cover~$X$.\nSo there are $n \\in {\\mathbb{Z}}_{> 0}$ and $g_1, g_2, \\ldots, g_n \\in G$\nsuch that the sets $g_1 U, g_2 U, \\ldots, g_n U$ cover~$X$.\nThe function $x \\mapsto \\sum_{k = 1}^n f (g_k^{-1} x)$\nis strictly positive on~$X$.\nUsing this fact at the first step, pure infiniteness of~$D$\nat the second last step,\nand (\\ref{Eq_0201_InD}) and~(\\ref{Eq_0201_Slice})\nat the last step,\nwe get\n\\begin{align*}\n1_A\n& \\precsim_{C (X)} \\sum_{k = 1}^n \\alpha_{g_k} (f \\otimes 1_D)\n\\\\\n& \\precsim_{C (X)} {\\mathrm{diag}} \\bigl( \\alpha_{g_1} (f \\otimes 1_D), \\,\n \\alpha_{g_2} (f \\otimes 1_D), \\, \\ldots, \\,\n \\alpha_{g_n} (f \\otimes 1_D) \\bigr)\n\\\\\n& = {\\mathrm{diag}} \\bigl( u_{g_1} (f \\otimes 1_D) u_{g_1}^*, \\,\n u_{g_2} (f \\otimes 1_D) u_{g_2}^*, \\, \\ldots, \\,\n u_{g_n} (f \\otimes 1_D) u_{g_n}^* \\bigr)\n\\\\\n& \\sim_{A} {\\mathrm{diag}} \\bigl( f \\otimes 1_D, \\,\n f \\otimes 1_D, \\, \\ldots, \\,\n f \\otimes 1_D \\bigr)\n = f \\otimes 1_{M_n (D)}\n \\precsim_{C (X, D)} f \\otimes 1_D\n \\precsim_{A} a.\n\\end{align*}\nThis completes the proof.\n\\end{proof}\n\nAs promised,\nwe now give a result on large subalgebras when $G = {\\mathbb{Z}}$.\nWe use two lemmas,\nthe first of which has a similar proof to that of\nTheorem~\\ref{T_0101_PICase}.\n\n\\begin{lem}\\label{L_0101_SubalgPI}\nLet $X$ be a compact metric space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a purely infinite simple unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nLet $Y \\subset X$ be a compact subset\nsuch that $h^n (Y) \\cap Y = \\varnothing$\nfor all $n \\in {\\mathbb{Z}} \\setminus \\{ 0 \\}$.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis purely infinite and simple.\n\\end{lem}\n\n\\begin{proof}\nFor convenience of notation,\nset $B = C^* \\big( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\big)_Y$.\nWe also freely identify $C (X, D)$ with $C (X) \\otimes D$.\nWe prove the result by showing that $1_A \\precsim a$\nfor all $a \\in A_{+} \\setminus \\{ 0 \\}$.\nBy Lemma~\\ref{L_4X21_79}\nand Lemma~\\ref{L_4X21_AYStruct},\nwe can assume that $a \\in C (X, D)_{+} \\setminus \\{ 0 \\}$.\nUsing Kirchberg's Slice Lemma\nand pure infiniteness of~$D$\nas in the proof of Theorem~\\ref{T_0101_PICase},\nthere is $f \\in C (X)_{+} \\setminus \\{ 0 \\}$\nsuch that $f \\otimes 1_D \\precsim_{C (X, D)} a$.\n\nWe now claim that for every $x \\in X$ there is $l_x \\in C (X)_{+}$\nsuch that $l_x \\otimes 1_D \\precsim_{B} f \\otimes 1_D$\nand $l_x (x) > \\frac{1}{2}$.\nGiven this,\nthe proof is finished as in the last computation\nin the proof of\nTheorem~\\ref{T_0101_PICase}.\n\nTo prove the claim,\nset $U = \\bigl\\{ x \\in X \\colon f (x) > \\frac{1}{2} \\bigr\\}$.\nIf $x \\in U$,\ntake $l_x = f$.\n\nSuppose next that $x \\not\\in U$\nand $x \\not\\in \\bigcup_{m = 0}^{\\infty} h^m (Y)$.\nBy minimality of~$h$,\nthere is $n \\in {\\mathbb{Z}}_{> 0}$ such that $x \\in h^{n} (U)$.\nChoose $r \\in C (X)_{+}$ such that $r (y) = 0$\nfor all $y \\in \\bigcup_{m = 0}^{n - 1} h^m (Y)$\nand $r (x) = 1$.\nThen, letting $u \\in C^* \\big( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\big)$\nbe the standard unitary (as in Notation~\\ref{ntn1-181221D}),\nwe have $(r \\otimes 1_D) u^n \\in B$\nby Lemma~\\ref{L_4X21_AYStruct}.\nTaking $l_x = r^2 \\cdot (f \\circ h^{- n})$,\nwe have $l_x (x) = f ( h^{- n} (x)) > \\frac{1}{2}$.\nAlso, $(r u^n) (f \\otimes 1_D) (r u^n)^* = l_x \\otimes 1_D$,\nso $l_x \\otimes 1_D \\precsim_{B} f \\otimes 1_D$.\n\nFinally, suppose $x \\not\\in U$\nand $x \\in \\bigcup_{m = 0}^{\\infty} h^m (Y)$.\nThen $x \\not\\in \\bigcup_{m = 1}^{\\infty} h^{- m} (Y)$.\nBy minimality of~$h$,\nthere is $n \\in {\\mathbb{Z}}_{> 0}$ such that $x \\in h^{- n} (U)$.\nChoose $r \\in C (X)_{+}$ such that $r (y) = 0$\nfor all $y \\in \\bigcup_{m = 1}^{n} h^{- m} (Y)$\nand $r (x) = 1$.\nThen $(r \\otimes 1_D) u^{- n} \\in B$\nby Lemma~\\ref{L_4X21_AYStruct}.\nTaking $l_x = r^2 \\cdot (f \\circ h^{n})$,\nwe have $l_x (x) = f ( h^{n} (x)) > \\frac{1}{2}$.\nAlso, $(r u^{- n}) (f \\otimes 1_D) (r u^{- n})^* = l_x \\otimes 1_D$,\nso $l_x \\otimes 1_D \\precsim_{B} f \\otimes 1_D$.\nThis completes the proof.\n\\end{proof}\n\n\\begin{lem}\\label{L_9Z01_KeepLarge}\nLet $X$ be a compact Hausdorff space,\nlet $G$ be a discrete group,\nand let $D$ be a C*-algebra.\nLet $(g, x) \\mapsto g x$ be an action of $G$ on~$X$,\nand let $\\alpha \\colon G \\to {\\mathrm{Aut}} (C (X, D))$\nbe an action of $G$ on $C (X, D)$\nwhich lies over $(g, x) \\mapsto g x$\nin the sense of Definition~\\ref{D_4Y12_OverX}.\nThen for every\n$b \\in C^*_{\\mathrm{r}} \\bigl( G, \\, C (X, D), \\, \\alpha \\bigr)$\nand every $\\varepsilon > 0$,\nthere is a finite set $T \\S G$ and a nonempty open set $V \\S X$\nsuch that,\nwhenever $x \\in V$\nand $f \\in C (X)$ satisfy $0 \\leq f \\leq 1$\nand $f (g x) = 1$ for all $g \\in T$,\nthen $\\| f b \\| > \\| b \\| - \\varepsilon$.\n\\end{lem}\n\nIf $D$ is not unital, then the product $f b$ is realized via\nthe inclusion of $C (X)$ in\n$M \\bigl( C^* \\bigl( G, \\, C (X, D), \\, \\alpha \\bigr) \\bigr)$.\n\n\\begin{proof}[Proof of Lemma~\\ref{L_9Z01_KeepLarge}]\nLet $(g, x) \\mapsto \\alpha_{g, x}$ be as in Definition~\\ref{D_4Y12_OverX}.\n\nFix a faithful representation $\\rho$ of $D$ on a Hilbert space~$H$.\nThen for every $x \\in X$ there is a representation\n$\\rho \\circ {\\mathrm{ev}}_x \\colon C (X, D) \\to L (H)$.\nLet $(v_x, \\pi_x)$ be the corresponding regular covariant\nrepresentation of $\\bigl( G, \\, C (X, D), \\, \\alpha \\bigr)$\non $l^2 (G, H)$,\nand let\n$\\sigma_x \\colon C^* \\bigl( G, \\, C (X, D), \\, \\alpha \\bigr) \\to l^2 (G, H)$\nbe its integrated form.\nWe identify $l^2 (G, H)$ with $l^2 (G) \\otimes H$,\nand we write $\\delta_g \\in l^2 (G)$ for the standard basis vector\ncorresponding to $g \\in G$.\nFor later use, we recall that if $S_1, S_2 \\S G$ are finite sets,\n$d \\in C (X, D) [G]$\nis given as $d = \\sum_{g \\in S_1} d_g u_g$\nwith $d_g \\in C (X, D)$ for $g \\in S_1$,\nand $\\xi \\in l^2 (G) \\otimes H$ has the form\n$\\xi = \\sum_{h \\in S_2} \\delta_h \\otimes \\xi_h$\nwith $\\xi_h \\in H$ for $h \\in S_2$,\nthen\n\\begin{equation}\\label{Eq_9X01_IndRepFormula}\n\\sigma_x (d) \\xi\n = \\sum_{g \\in S_1} \\sum_{h \\in S_2}\n \\delta_{g h} \\otimes\n \\rho \\bigl( \\alpha_{h^{-1} g^{-1}, x} (d_g (g h x)) \\bigr) \\xi_h.\n\\end{equation}\n\nThe representation $\\bigoplus_{x \\in X} \\rho \\circ {\\mathrm{ev}}_x$\nis a faithful representation of $C (X, D)$,\nso $\\bigoplus_{x \\in X} \\sigma_x$\nis a faithful representation of\n$C^*_{\\mathrm{r}} \\bigl( G, \\, C (X, D), \\, \\alpha \\bigr)$.\nTherefore there exists $x_0 \\in X$ such that\n$\\| \\sigma_{x_0} (b) \\| > \\| b \\| - \\frac{\\varepsilon}{5}$.\nChoose $c \\in C (X, D) [G]$ such that\n$\\| b - c \\| < \\frac{\\varepsilon}{5}$.\nIn particular,\n$\\| \\sigma_{x_0} (c) \\| > \\| b \\| - \\frac{2 \\varepsilon}{5}$.\nChoose $\\xi \\in l^2 (G, H)$ with finite support such that\n$\\| \\xi \\| = 1$ and\n$\\| \\sigma_{x_0} (c) \\xi \\| > \\| b \\| - \\frac{3 \\varepsilon}{5}$.\nWrite $c = \\sum_{g \\in S_1} c_g u_g$\nand $\\xi = \\sum_{h \\in S_2} \\delta_h \\otimes \\xi_h$\nwith $S_1, S_2 \\S G$ finite,\n$c_g \\in C (X, D)$ for $g \\in S_1$,\nand $\\xi_h \\in H$ for $h \\in S_2$.\nUse Lemma~\\ref{L_9Y30_MixCont} to choose\nan open set $V \\S X$ such that $x_0 \\in V$\nand such that for all $x \\in V$, $g \\in S_1$, and $h \\in S_2$, we have\n\\begin{equation}\\label{Eq_9Z01_Choice}\n\\bigl\\| \\alpha_{h^{-1} g^{-1}, x} (c_g (g h x))\n - \\alpha_{h^{-1} g^{-1}, x_0} (c_g (g h x_0) ) \\bigr\\|\n < \\frac{\\varepsilon}{5 {\\mathrm{card}} (S_1) {\\mathrm{card}} (S_2) + 1}.\n\\end{equation}\nSet\n\\[\nT = \\bigl\\{ g h \\colon {\\mbox{$g \\in S_1$ and $h \\in S_2$}} \\bigr\\}.\n\\]\nNow let $x \\in V$,\nand suppose $f \\in C (X)$ satisfies $0 \\leq f \\leq 1$\nand $f (g x) = 1$ for all $g \\in T$.\nThe condition on~$f$\nand the formula~(\\ref{Eq_9X01_IndRepFormula})\nimply that $\\sigma_x (f c) \\xi = \\sigma_x (c) \\xi$.\nApplying~(\\ref{Eq_9X01_IndRepFormula}) again,\nand using~(\\ref{Eq_9Z01_Choice}) at the second step,\nwe get\n\\begin{equation}\\label{Eq_9Z01_xx0}\n\\| \\sigma_x (c) \\xi - \\sigma_{x_0} (c) \\xi \\|\n \\leq \\sum_{g \\in S_1} \\sum_{h \\in S_2}\n \\bigl\\| \\alpha_{h^{-1} g^{-1}, x} (c_g (g h x))\n - \\alpha_{h^{-1} g^{-1}, x_0} (c_g (g h x_0) ) \\bigr\\|\n < \\frac{\\varepsilon}{5}.\n\\end{equation}\nTherefore, using $\\| f \\| \\leq 1$ at the second step\nand $\\| b - c \\| < \\frac{\\varepsilon}{5}$ and the second and fifth steps,\nas well as (\\ref{Eq_9Z01_xx0}) at the fourth step,\n\\begin{align*}\n\\| f b \\|\n& \\geq \\| \\sigma_x (f b) \\|\n > \\| \\sigma_x (f c) \\xi \\| - \\frac{\\varepsilon}{5}\n = \\| \\sigma_x (c) \\xi \\| - \\frac{\\varepsilon}{5}\n > \\| \\sigma_{x_0} (c) \\xi \\| - \\frac{2 \\varepsilon}{5}\n\\\\\n& > \\| \\sigma_{x_0} (b) \\xi \\| - \\frac{3 \\varepsilon}{5}\n > \\| \\sigma_{x_0} (b) \\| - \\frac{4 \\varepsilon}{5}\n > \\| \\sigma_{x_0} (b) \\| - \\frac{4 \\varepsilon}{5}\n > \\| b \\| - \\varepsilon,\n\\end{align*}\nas desired.\n\\end{proof}\n\n\\begin{prp}\\label{T_9921_PICase}\nLet $X$ be a compact metric space,\nlet $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $Y \\subset X$ be a compact subset\nsuch that $h^{n} (Y) \\cap Y = \\varnothing$\nfor all $n \\in {\\mathbb{Z}} \\setminus \\set{0}$,\nlet $D$ be a purely infinite simple unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nThen\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis a large subalgebra of\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nof crossed product type.\n\\end{prp}\n\n\\begin{proof}\nWe verify directly the conditions of Definition~4.9 of~\\cite{PhLg}.\nWe follow Notation~\\ref{N_4Y12_CPN}.\nSet\n\\[\nA = C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr),\n\\qquad\nB = C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y,\n\\]\n\\[\nC = C (X, D),\n\\qquad {\\mbox{and}} \\qquad\nG = \\{ u \\}.\n\\]\nAll parts of condition~(1)\nin Definition~4.9 of~\\cite{PhLg}\nare obvious.\n\nWe verify condition~(2) there.\nLet $m \\in {\\mathbb{Z}}_{> 0}$,\nlet $a_1, a_2, \\ldots, a_m \\in A$,\nlet $\\varepsilon > 0$,\nlet $b_1 \\in A_{+}$ satisfy $\\| b_1 \\| = 1$,\nand let $b_2 \\in B_{+} \\setminus \\{ 0 ]$.\nWe will part of the verification of\ncondition~(3)\nin Proposition 4.11 of~\\cite{PhLg}\nin the proof of Theorem~\\ref{T_5418_AYLg}.\nChoose $c_1, c_2, \\ldots, c_m \\in C (X, D) [{\\mathbb{Z}}]$\nsuch that $\\| c_j - a_j \\| < \\varepsilon$\nfor $j = 1, 2, \\ldots, m$.\n(This estimate is condition~(2b).)\nChoose $N \\in {\\mathbb{Z}}_{> 0}$\nsuch that for $j = 1, 2, \\ldots, m$\nthere are $c_{j, l} \\in C (X, D)$\nfor $l = -N, \\, - N + 1, \\, \\ldots, \\, N - 1, \\, N$\nwith\n\\[\nc_j = \\sum_{l = -N}^N c_{j, l} u^l.\n\\]\nApply Lemma~\\ref{L_9Z01_KeepLarge},\ngetting a finite set $T \\S {\\mathbb{Z}}$ and a nonempty open set $U \\S X$\nsuch that,\nwhenever $x \\in U$\nand $f \\in C (X)$ satisfy $0 \\leq f \\leq 1$\nand $f (h^k (x)) = 1$ for all $k \\in T$,\nthen $\\| f b_1^{1 \/ 2} \\| > 1 - \\frac{\\varepsilon}{2}$.\nDefine\n\\[\nS = \\bigl\\{ n - k \\colon\n {\\mbox{$k \\in T$ and $n \\in [- N, N] \\cap {\\mathbb{Z}}$}} \\bigr\\},\n\\]\nwhich is a finite subset of~${\\mathbb{Z}}$.\nThen $\\bigcup_{n \\in S} h^n (Y)$ is nowhere dense in~$X$,\nso there exists $x_0 \\in W$\nsuch that $x_0 \\not\\in \\bigcup_{n \\in S} h^n (Y)$.\nThis choice implies that\n\\[\n\\{ h^k (x_0) \\colon k \\in T \\} \\cap \\bigcup_{n = - N}^{N} h^n (Y)\n = \\varnothing,\n\\]\nso there is $g \\in C (X)$\nsuch that $0 \\leq g \\leq 1$,\n$g (y) = 1$ for all $y \\in \\bigcup_{n = - N}^{N} h^n (Y)$,\nand $g (h^k (x)) = 0$ for all $k \\in T$.\n\nCondition~(1a) in Definition~4.9 of~\\cite{PhLg}\nholds by construction.\nThe proof that $(1 - g) c_j \\in B$ for $j = 1, 2, \\ldots, m$\n(condition~(2c) in Definition~4.9 of~\\cite{PhLg})\nis the same as at the end of the proof of Theorem~\\ref{T_5418_AYLg}:\nfollow\nthe corresponding part of the proof of Theorem~7.10 of~\\cite{PhLg},\nexcept using \\Lem{L_4X21_AYStruct}\nin place of Proposition~7.5 of~\\cite{PhLg}.\nLemma~\\ref{L_0101_SubalgPI}\nimplies $1_A \\precsim_A b_1$ and $1_A \\precsim_B b_2$,\nso the requirements $g \\precsim_A b_1$ and $g \\precsim_B b_2$\n(condition~(2d) in Definition~4.9 of~\\cite{PhLg})\nfollow immediately.\nFor condition~(2e) in Definition~4.9 of~\\cite{PhLg},\nwe estimate, using $1 - g (h^k (x_0)) = 1$\nfor all $k \\in T$ at the second step:\n\\[\n\\bigl\\| (1 - g) b_1 (1 - g) \\bigr\\|\n = \\| (1 - g) b_1^{1 \/ 2} \\|^2\n > \\bigl( 1 - \\frac{\\varepsilon}{2} \\bigr)^2\n > \\varepsilon.\n\\]\nThis completes the proof.\n\\end{proof}\n\n\\section{Recursive structure for orbit breaking\n subalgebras}\\label{RecStruct}\n\nIn this section, under appropriate conditions\nwe will show that the orbit breaking subalgebras\nof Definition~\\ref{D_4Y12_OrbSubalg}\nhave a tractable recursive structure,\nas subalgebras of certain homogeneous algebras\ntensored with $D$. We start by introducing the\nformalism for a generalization of the recursive\nsubhomogeneous algebras introduced in\n\\cite{PhRsha1} that were crucial for the analysis in\n\\cite{QLinPhDiff} and \\cite{HLinPh}.\n\n\\begin{dfn}\\label{pullback}\nLet $A$, $B$, and $C$ be C*-algebra{s}, and let $\\varphi \\colon\nA \\to C$ and $\\psi \\colon B \\to C$ be homomorphisms.\nThen the associated {\\emph{pullback C*-algebra{}}} $A\n\\oplus_{C, \\varphi, \\psi} B$ is defined by\n\\[\nA \\oplus_{C, \\varphi, \\psi} B\n = \\bset{ (a, b) \\in A \\oplus B \\colon \\varphi (a) = \\psi (b) }.\n\\]\n\\end{dfn}\n\nWe frequently write $A \\oplus_{C} B$ when the maps $\\varphi$\nand $\\psi$ are understood.\nIf $C = 0$ we just get $A \\oplus B$.\n\n\\begin{dfn}\\label{RSHA}\nLet $D$ be a simple unital C*-algebra.\nThe class of\n{\\emph{recursive subhomogeneous algebras over $D$}} is\nthe smallest class $\\mathcal{R}$ of C*-algebra{s} that is\nclosed under isomorphism and such that:\n\\begin{enumerate}\n\\item\\label{9223_RSHA_1}\nIf $X$ is a compact Hausdorff space and $n \\geq 1$,\nthen $C (X, M_{n} (D)) \\in \\mathcal{R}$.\n\\item\\label{9223_RSHA_2}\nIf $B \\in \\mathcal{R}$, $X$ is compact Hausdorff,\n$n \\geq 1$, $X^{(0)} \\subset X$ is closed (possibly empty),\n$\\varphi \\colon B \\to C (X^{(0)}, \\, M_{n} (D))$\nis any unital homomorphism\n(the zero homomorphism{} if $X^{(0)}$ is empty), and\n$\\rho \\colon C (X, M_{n} (D)) \\to C (X^{(0)}, \\, M_{n} (D))$ is the\nrestriction homomorphism, then the pullback\n\\begin{align*}\n& B \\oplus_{C (X^{(0)}, \\, M_{n} (D)), \\, \\varphi, \\, \\rho)} C (X, M_{n} (D))\n\\\\\n& \\hspace*{3em} {\\mbox{}}\n = \\bset{ (b, f) \\in B \\oplus C (X, M_{n} (D)) \\colon\n \\varphi (b) = f |_{(X^{(0)}} }\n\\end{align*}\nis in $\\mathcal{R}$.\n\\end{enumerate}\n\\end{dfn}\n\nTaking $D = {\\mathbb{C}}$ in this definition gives the usual definition\nfor the class of recursive subhomogeneous algebras.\n(See \\cite{PhRsha1}.)\nThis definition makes sense for any unital C*-algebra~$D$,\nbut it is not clear whether it is appropriate in this generality.\n\n\\begin{dfn}\\label{RSHAMachinery}\nWe adopt the following standard notation for recursive\nsubhomogeneous algebras over $D$.\nThe definition implies\nthat if $R$ is any recursive subhomogeneous algebra over\n$D$, then $R$ may be written in the form\n\\[\nR \\cong\n \\left[ \\cdots \\left[ \\left[ C_{0} \\oplus_{C_{1}^{(0)}, \\varphi_1, \\rho_1}\n C_{1} \\right]\n \\oplus_{C_{2}^{(0)}, \\varphi_2, \\rho_2} \\right] \\cdots \\right]\n \\oplus_{C_{l}^{(0)}, \\varphi_l, \\rho_l} C_{l},\n\\]\nwith $C_{k} = C (X_{k}, \\, M_{n (k)} (D))$ for compact Hausdorff\nspaces $X_{k}$ and positive integers $n (k)$, and with\n$C_{k}^{(0)} = C \\bigl( X_{k}^{(0)}, \\, M_{n (k)} (D) \\bigr)$\nfor compact\nsubsets $X_{k}^{(0)} \\subset X_{k}$ (possibly empty),\nwhere the maps $\\rho_{k} \\colon C_{k} \\to C_{k}^{(0)}$ are\nalways the restriction maps.\nAn expression of this type for\n$R$ will be referred to as a {\\emph{decomposition}} of~$R$,\nand the notation that appears here will be referred to as the\n{\\emph{standard notation for a decomposition}}.\nWe associate the following objects to this decomposition.\n\\begin{enumerate}\n\\item\\label{9223_RSHAMachinery_1}\nIts {\\emph{length}} $l$.\n\\item\\label{9223_RSHAMachinery_2}\nThe {\\emph{$k$-th stage algebra}}\n\\[\nR^{(k)}\n = \\left[ \\cdots \\left[ \\left[ C_{0} \\oplus_{C_{1}^{(0)}}\n C_{1} \\right] \\oplus_{C_{2}^{(0)}} C_{2} \\right] \\cdots \\right]\n \\oplus_{C_{k}^{(0)}} C_{k}.\n\\]\n\\item\\label{9223_RSHAMachinery_3}\nIts {\\emph{base spaces}} $X_{0},\\ldots,X_{l}$ and {\\emph{total space}}\n$X = \\coprod_{k = 0}^{l} X_{k}$.\n\\item\\label{9223_RSHAMachinery_4}\nIts {\\emph{matrix sizes}} $n (0),\\ldots,n (l)$ and {\\emph{matrix\nsize function}} $m \\colon X \\to {\\mathbb{Z}}_{\\geq 0}$ defined by $m(x) = n (k)$\nwhen $x \\in X_{k}$.\n(This is called the {\\emph{matrix size of $R$ at~$x$}}.)\n\\item\\label{9223_RSHAMachinery_5}\nIts {\\emph{minimum matrix size}} $\\min_{k} n (k)$ and\n{\\emph{maximum matrix size}} $\\max_{k} n (k)$.\n\\item\\label{9223_RSHAMachinery_6}\nIts {\\emph{topological dimension}}\n$\\dim (X) = \\max_{k} \\dim (X_{k})$\nand {\\emph{topological dimension function}}\n$d \\colon X \\to {\\mathbb{Z}}_{\\geq 0}$, defined by $d(x) = \\dim (X_{k})$\nfor $x \\in X_{k}$.\n(This is called the {\\emph{topological dimension of $R$ at $x$}}.)\n\\item\\label{9223_RSHAMachinery_7}\nIts {\\emph{standard representation}}\n$\\sigma\n = \\sigma_{R} \\colon R\n \\to \\bigoplus_{k = 0}^{l} C (X_{k}, \\, M_{n (k)} (D))$,\ndefined by\nforgetting the restriction to a subalgebra in each of the fibered\nproducts in the decomposition.\n\\item\\label{9223_RSHAMachinery_8}\nThe associated {\\emph{evaluation maps}}\n${\\mathrm{ev}}_{x} \\colon R \\to M_{n (k)} (D)$,\ndefined to be the restriction of the usual evaluation\nmap on $\\bigoplus_{k = 0}^{l} C (X_{k}, \\, M_{n (k)} (D))$ to $R$,\nwhen $R$\nis identified with a subalgebra of this algebra through the standard\nrepresentation $\\sigma_{R}$.\n\\end{enumerate}\n\\end{dfn}\n\nA decomposition of an algebra $R$ as a\nrecursive subhomogeneous algebra over $D$ may be highly\nnonunique, which in the case $D = {\\mathbb{C}}$ is clear from examples\nin~\\cite{PhRsha1}.\nMoreover, the matrix sizes are not uniquely\ndetermined even if all the other data has already been chosen,\nwhich is easily seen through\nthe example of the $2^{\\infty}$~UHF algebra.\n\n\\begin{ntn}\\label{StronglySelfAbs}\nThroughout, we let $\\mathcal{E}$ denote a separable, unital C*-algebra{}\nthat is strongly selfabsorbing in the sense of \\cite{TomsWinter2}:\nthere\nexists an isomorphism $\\mathcal{E} \\cong \\mathcal{E} \\otimes \\mathcal{E}$ that is unitarily\nequivalent to $a \\mapsto a \\otimes 1_{\\mathcal{E}}$.\n\\end{ntn}\n\nExamples include the Cuntz algebras\n$\\mathcal{O}_{2}$ and $\\mathcal{O}_{\\infty}$,\nand the Jiang-Su algebra ${\\mathcal{Z}}$,\nwhich is a simple, separable, unital, infinite dimensional,\nstrongly selfabsorbing, nuclear C*-algebra{} having the same\nElliott invariant\nas the complex numbers ${\\mathbb{C}}$.\n(See \\cite{JiangSu} for its construction).\n\n\\begin{dfn}\\label{DStableDefn}\nAdopt Notation~\\ref{StronglySelfAbs}.\nA separable C*-algebra{} $D$ is called {\\emph{$\\mathcal{E}$-stable}}\nif there is\nan isomorphism $D \\otimes \\mathcal{E} \\cong D$.\n\\end{dfn}\n\nIt is clear that if $D$ is $\\mathcal{E}$-stable, then so are $M_{n} (D)$\nand $C (X, D)$.\nWith appropriate assumptions on the algebra $D$, we can give some\nresults about the $\\mathcal{E}$-stability of recursive subhomogeneous algebras\nover $C$.\n\n\\begin{lem}\\label{DStablePullback}\nAdopt Notation~\\ref{StronglySelfAbs}.\nLet $A$, $B$, and $C$ be separable unital C*-algebra{s}\n(allowing $C = 0$),\nand let $\\varphi \\colon A \\to C$ and $\\psi \\colon B \\to C$\nbe unital homomorphisms.\nLet $P = A \\oplus_{C, \\varphi, \\psi} B$ be the associated pullback\n(Definition~\\ref{pullback}).\nIf $\\psi$\nis surjective and both $A$ and $B$ are $\\mathcal{E}$-stable, then $P$ is\n$\\mathcal{E}$-stable.\n\\end{lem}\n\n\\begin{proof}\nDefine\n$\\pi \\colon P \\to A$ by $\\pi (a, b) = a$.\nThere is an isomorphism\n$\\iota \\colon \\mathrm{Ker} (\\psi) \\to \\mathrm{Ker} (\\pi)$\ngiven by $\\iota (b) = (0, b)$ for $b \\in \\mathrm{Ker} (\\psi)$.\nMoreover, surjectivity of $\\psi$ is easily seen to imply\nsurjectivity of $\\pi$.\nWe obtain an exact\nsequence\n\\[\n0 \\longrightarrow \\mathrm{Ker} (\\psi)\n \\stackrel{\\iota}{\\longrightarrow} P\n \\stackrel{\\pi}{\\longrightarrow} A\n \\longrightarrow 0.\n\\]\nNow $\\mathrm{Ker} (\\varphi)$ is $\\mathcal{E}$-stable by Corollary 3.1 of\n\\cite{TomsWinter2}.\nSince $B$ is also $\\mathcal{E}$-stable, Theorem 4.3 of\n\\cite{TomsWinter2} implies that $P$ is $\\mathcal{E}$-stable.\n\\end{proof}\n\n\n\\begin{prp}\\label{DStableRSHA}\nAdopt Notation~\\ref{StronglySelfAbs}.\nLet $D$ be a simple separable $\\mathcal{E}$-stable C*-algebra,\nand let $R$ be a recursive subhomogeneous algebra over $D$.\nThen $R$ is $\\mathcal{E}$-stable.\n\\end{prp}\n\n\\begin{proof}\nWe proceed by induction on the length of a decomposition for $R$ as\na recursive subhomogeneous algebra over $D$.\nThe base case is when\n$R = C (X, M_{n} (D))$, and this is $\\mathcal{E}$-stable whenever $D$ is\n$\\mathcal{E}$-stable.\nFor the inductive step, we may assume that\nthere are an $\\mathcal{E}$-stable unital C*-algebra~$R'$,\n$n \\in {\\mathbb{Z}}_{> 0}$, a compact Hausdorff space~$X$,\na closed subset $X^{(0)} \\subset X$,\nand a unital homomorphism\n$\\varphi \\colon R' \\to C \\bigl( X^{(0)}, \\, M_{n} (D) \\bigr)$\nsuch that\n\\[\nR = \\bset{ (a, f) \\in R' \\oplus C (X, M_{n} (D))\n \\colon \\varphi (a) = f |_{X^{(0)}} }\n\\]\nSince $f \\mapsto f |_{X^{(0)}}$ is surjective and both\n$R'$ and $C (X, M_{n} (D))$\nare $\\mathcal{E}$-stable, it follows from Lemma~\\ref{DStablePullback}\nthat $R$ is $\\mathcal{E}$-stable,\nwhich completes the induction.\n\\end{proof}\n\n\\begin{prp}\\label{AyDirLim}\nLet $X$ be a compact Hausdorff space, let $h \\colon X \\to X$ be a homeomorphism,\nlet $D$ be a unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nLet $Y \\subset X$ be closed, and let\n$Y_{1} \\supset Y_{2} \\supset \\cdots$\nbe closed subsets of $X$ such that $\\bigcap_{n = 1}^{\\infty}\nY_{n} = Y$.\nThen\n\\[\nC^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y\n = {\\overline{ {\\ts{ {\\ds{\\bigcup}}_{n = 1}^{\\infty} }} A_{Y_{n}} }}\n \\cong \\varinjlim A_{Y_{n}}.\n\\]\n\\end{prp}\n\n\\begin{proof}\nThe proof is easy, and is omitted.\n\\end{proof}\n\nThe results in the remainder of this section are\nmostly generalizations of ones in Section 1 of\n\\cite{QLinPhDiff}.\nWe follow a slightly more modern development,\nadapted from Section 11.3 of \\cite{PhNotes}, since \\cite{QLinPhDiff}\nhas not been published.\n(The proofs in \\cite{PhNotes} are mostly\ntaken from \\cite{QLinPhDiff}, with some changes in notational\nconventions.)\nMost proofs go through with changes only to the notation,\nsuch as replacing the action of a minimal homeomorphism\n$h$ with the action $\\alpha$ on $C (X, D)$.\nThe biggest technical\ndifferences are in the proof of Lemma~\\ref{L_4X21_AYStruct}.\nWe begin with the well-known Rokhlin tower construction.\n\n\\begin{ntn}\\label{ModifiedRokhlinTower}\nLet $X$ be a compact Hausdorff space{} and let $h \\colon X \\to X$ be a minimal homeomorphism.\nLet $Y \\subset X$ be a closed set with ${\\mathrm{int}} (Y) \\neq \\varnothing$.\nFor $y \\in Y$, define\n$r (y) = \\inf \\bigl( \\set{m \\geq 1 \\colon h^{m} (y) \\in Y} \\bigr)$.\nThen $\\sup_{y \\in Y} r (y) < \\infty$\n(see Lemma \\ref{RokhlinTowerProps}(\\ref{9Z26_new})),\nso there are finitely many\ndistinct values $n (0) < n (1) < \\cdots < n (l)$ in the range of~$r$.\nFor $k = 0, 1, \\ldots l$, set\n\\[\nY_{k} = \\overline{ \\set{ y \\in Y \\colon r (y) = n (k) } }\n\\qquad {\\mbox{and}} \\qquad\nY_{k}^{\\circ}\n = {\\mathrm{int}} \\bigl( \\bset{ y \\in Y \\colon r (y) = n (k) } \\bigr).\n\\]\n\\end{ntn}\n\nWe warn that, while\n$Y_{k}^{\\circ} \\subset {\\mathrm{int}} (Y_k)$,\nthe inclusion may be proper.\n\n\\begin{lem}\\label{RokhlinTowerProps}\nLet $Y \\subset X$ be a closed set with ${\\mathrm{int}} (Y) \\neq \\varnothing$,\nand adopt Notation \\ref{ModifiedRokhlinTower}.\nThen:\n\\begin{enumerate}\n\\item\\label{9Z26_new}\n$\\sup_{y \\in Y} r (y) < \\infty$.\n\\item\\label{Item_RokhlinTowerProps_1}\nThe sets $h^{j} (Y_{k}^{\\circ})$ are pairwise disjoint for $0\n\\leq k \\leq l$ and $1 \\leq j \\leq n (k)$.\n\\item\\label{Item_RokhlinTowerProps_2}\n$\\bigcup_{k = 0}^{l} Y_{k} = Y$.\n\\item\\label{Item_RokhlinTowerProps_3}\n$\\bigcup_{k = 0}^{l} \\bigcup_{j= 0}^{n (k)-1} h^{j} (Y_{k}) = X$.\n\\end{enumerate}\n\\end{lem}\n\n\\begin{proof}\nThis is Lemma 11.3.5 of~\\cite{PhNotes}.\n\\end{proof}\n\nWe first consider the specific situation of the structure of\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nwhen $X$ is the Cantor set.\nThis is needed later.\n\n\\begin{lem}\\label{CantorSetCase}\nAdopt Notation \\ref{ModifiedRokhlinTower}, let $X$ be the Cantor\nset, and let $Y \\subset X$ be a nonempty compact open subset.\nThen\n\\[\nC^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y\n \\cong \\bigoplus_{k = 0}^{l} C (Y_{k}, \\, M_{n (k)} (D)).\n\\]\n\\end{lem}\n\n\\begin{proof}\nThis is a straightforward adaption of the proof of Lemma 11.2.20\nof~\\cite{PhNotes}.\n\\end{proof}\n\nIf we set $B_{k} = C (Y_{k}, M_{n (k)})$, then $B_{k}$ is an\nAF algebra for each~$k$\n(since $Y_{k}$ is totally disconnected), and hence\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis a direct sum of what might be called ``AF $D$-algebras''.\n\n\\begin{prp}\\label{BanachSpaceDirSum}\nLet $X$ be a compact Hausdorff space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a unital C*-algebra,\nlet $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$,\nlet $Y \\subset X$ be a closed set with ${\\mathrm{int}} (Y) \\neq \\varnothing$,\nand adopt Notation \\ref{ModifiedRokhlinTower}.\nFollowing the notation of Lemma~\\ref{L_4X21_AYStruct}, there is\n$N \\in {\\mathbb{Z}}_{> 0}$ such that $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nhas the Banach space direct sum decomposition\n\\[\nC^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y\n = \\bigoplus_{n = - N}^{N} C_{0} (X \\setminus Y_{n}, \\, D) u^{n}.\n\\]\n\\end{prp}\n\n\\begin{proof}\nThe proof is the same as that of Corollary 11.3.7 of~\\cite{PhNotes}.\n\\end{proof}\n\n\\begin{prp}\\label{AyHom}\nLet $X$ be a compact Hausdorff space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a unital C*-algebra,\nlet $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$,\nlet $Y \\subset X$ be a closed set with ${\\mathrm{int}} (Y) \\neq \\varnothing$,\nand adopt Notation \\ref{ModifiedRokhlinTower}.\nDefine the unitary $s_{k} \\in C (Y_{k}, M_{n (k)} (D))$ by\n\\[\ns_{k} = \\left[ \\begin{array}{ccccccc}\n0 & 0 & \\cdots & \\cdots & 0 & 0 & 1 \\\\\n1 & 0 & \\cdots & \\cdots & 0 & 0 & 0 \\\\\n0 & 1 & \\cdots & \\cdots & 0 & 0 & 0 \\\\\n\\vdots & \\vdots & \\ddots & & \\vdots & \\vdots & \\vdots \\\\\n\\vdots & \\vdots & & \\ddots & \\vdots & \\vdots & \\vdots \\\\\n0 & 0 & \\cdots & \\cdots & 1 & 0 & 0 \\\\\n0 & 0 & \\cdots & \\cdots & 0 & 1 & 0 \\end{array} \\right].\n\\]\nFor $k = 0, 1, \\ldots, l$ there is a unique linear map\n\\[\n\\gamma_{k} \\colon C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y\n \\to C (Y_{k}, M_{n (k)} (D))\n\\]\nsuch that for $m = 0, 1, \\ldots, n (k) - 1$\nand $f \\in C_{0} (X \\setminus Z_{m}, D)$ we have:\n\\begin{enumerate}\n\\item\\label{9222_AyHom_1}\n$\\gamma_{k} (f u^{m}) = {\\mathrm{diag}} \\bigl( f |_{Y_{k}}, \\,\n \\alpha^{-1} (f) |_{Y_{k}}, \\,\n \\ldots, \\, \\alpha^{-[n (k)-1]} (f) |_{Y_{k}} \\bigr) \\cdot s_{k}^{m}$.\n\\item\\label{9222_AyHom_2}\n$\\gamma_{k} (u^{-m} f) = s_{k}^{-m} \\cdot\n {\\mathrm{diag}} \\bigl( f |_{Y_{k}}, \\, \\alpha^{-1} (f)\n |_{Y_{k}}, \\, \\ldots, \\,\n \\alpha^{-[n (k)-1]} (f) |_{Y_{k}} \\bigr)$.\n\\end{enumerate}\nMoreover, the map\n\\[\n\\gamma \\colon C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y\n \\to \\bigoplus_{k = 0}^{l} C (Y_{k}, M_{n (k)} (D))\n\\]\ngiven by $\\gamma (a) = (\\gamma_{0} (a), \\gamma_{1} (a), \\ldots, \\gamma_{l} (a))$\nis a\n$*$-homomorphism.\n\\end{prp}\n\n\\begin{proof}\nThe computations in this proof are analogous to those in\nthe proof of Proposition 11.3.9 of~\\cite{PhNotes}, using\nLemma~\\ref{L_4X21_AYStruct}\nin place of Proposition 11.3.6 of~\\cite{PhNotes}.\n\\end{proof}\n\n\\begin{lem}\\label{SubdiagProjnDecomp}\nLet $X$ be a compact Hausdorff space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a unital C*-algebra,\nlet $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$,\nand let $Y \\subset X$ be a closed set with ${\\mathrm{int}} (Y) \\neq \\varnothing$.\nAdopt Notation \\ref{ModifiedRokhlinTower}.\nIdentify $C (Y_{k}, \\, M_{n (k)} (D))$\nwith $M_{n (k)} (C (Y_{k}, D))$ in the\nobvious way.\nDefine maps\n\\[\nE_{k}^{(m)} \\colon C (Y_{k}, M_{n (k)} (D))\n\\to C (Y_{k}, M_{n (k)} (D))\n\\]\nby $E_{k}^{(m)} (b)_{m+j,j} = b_{m+j,j}$\nfor $1 \\leq j \\leq n (k) - m$ (if $m \\geq 0$)\nand for $-m + 1 \\leq j \\leq n (k)$ (if $m \\leq 0$),\nand $E_{k}^{(m)} (b)_{i,j} = 0$ for\nall other pairs $(i,j)$.\n(Thus, $E_{k}^{(m)}$ is the projection map\non the $m$-th subdiagonal.)\nWrite\n\\[\nG_{m} = \\bigoplus_{k = 0}^{l} E_{k}^{(m)} (C (Y_{k}, M_{n (k)} (D)).\n\\]\nThen:\n\\begin{enumerate}\n\\item\\label{SubdiagProjnDecomp_DirSum}\nThere is a Banach space direct sum decomposition\n\\[\n\\bigoplus_{k = 0}^{l} C (Y_{k}, M_{n (k)} (D)) =\n\\bigoplus_{m=-n (l)}^{n (l)} G_{m}.\n\\]\n\\item\\label{SubdiagProjnDecomp_gmIs}\nFor $k = 0, 1,\\ldots,l$, $m \\geq 0$,\n$f \\in C_{0} (X \\setminus Z_{m}, D)$,\nand $x \\in Y_{k}$,\nthe expression $\\gamma_{k} (f u^{m}) (x)$\nis given by the following matrix, in which the first nonzero\nentry is in row $m+1$:\n\\[\n\\gamma_{k} (f u^{m}) (x)\n = \\left[ \\begin{array}{ccccccc} 0 & 0 & \\cdots &\n\\cdots & \\cdots & \\cdots & 0\n\\\\\n\\vdots & \\vdots & & & & & \\vdots \\\\\n0 & 0 & & & & & \\vdots\n\\\\\n\\alpha^{-m} (f) (x) & 0 & & & & & \\vdots \\\\\n0 & \\alpha^{-[m+1]} (f) (x) & & & & & \\vdots\n\\\\\n\\vdots & & \\ddots & & &\n& \\vdots\n\\\\\n0 & \\cdots & \\cdots & \\alpha^{-[n (k)-1]} (f) (x) & 0\n& \\cdots & 0\n\\end{array} \\right].\n\\]\n\\item\\label{SubdiagProjnDecomp_3}\nFor $m \\geq 0$ and $f \\in C_{0} (X \\setminus Z_{m}, D)$, we\nhave\n\\[\n\\gamma_{k} (f u^{m}) \\in E_{k}^{(m)} (C (Y_{k}, M_{n (k)} (D))),\n\\qquad\n\\gamma (f u^{m}) \\in G_{m},\n\\]\n\\[\n\\gamma_{k} (u^{-m} f) \\in E_{k}^{(-m)} (C (Y_{k}, M_{n (k)} (D))),\n\\qquad {\\mbox{and}} \\qquad\n\\gamma (u^{-m} f) \\in G_{-m}.\n\\]\n\\item\\label{SubdiagProjnDecomp_Compatible}\nThe homomorphism $\\gamma$ is compatible with the\ndirect sum decomposition of Proposition~\\ref{BanachSpaceDirSum}\non its domain and the direct sum decomposition of part (1) on its\ncodomain.\n\\end{enumerate}\n\\end{lem}\n\n\\begin{proof}\nThe direct sum decomposition is essentially immediate from the\ndefinition of the maps $E_{k}^{(m)}$, while the other statements\nfollow from Proposition~\\ref{BanachSpaceDirSum} and some\nstraightforward matrix calculations as in the proof of Lemma 11.3.15\nof~\\cite{PhNotes}.\n\\end{proof}\n\n\\begin{cor}\\label{GammaInj}\nLet $X$ be a compact Hausdorff space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nLet $Y \\subset X$ be closed with ${\\mathrm{int}} (Y) \\neq \\varnothing$.\nThen\nthe homomorphism $\\gamma$ of Proposition~\\ref{BanachSpaceDirSum}\nis injective.\n\\end{cor}\n\n\\begin{proof}\nThe proof is the same as the proof of Lemma 11.3.17 of~\\cite{PhNotes}.\n\\end{proof}\n\n\\begin{lem}\\label{ByMemberCond}\nLet $X$ be a compact Hausdorff space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nLet $Y \\subset X$ be closed with ${\\mathrm{int}} (Y) \\neq \\varnothing$.\nLet\n\\[\nb = (b_{0}, b_{1}, \\ldots, b_{l})\n \\in \\bigoplus_{k = 0}^{l} C (Y_{k}, M_{n (k)} (D)).\n\\]\nThen\n$b \\in \\gamma \\bigl( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y \\bigr)$\nif and only if whenever\n\\begin{itemize}\n\\item\n$r > 0$,\n\\item\n$k,t_{1},\\ldots,t_{r} \\in \\set{0,\\ldots,l}$,\n\\item\n$n (t_{1}) + n (t_{2}) + \\cdots + n (t_{r}) = n (k)$,\n\\item\n$x \\in (Y_{k} \\setminus Y_{k}^{\\circ}) \\cap\nY_{t_{1}} \\cap h^{-n (t_{1})} (Y_{t_{2}}) \\cap \\cdots \\cap\nh^{-[n (t_{1}) + \\cdots + n (t_{r - 1}) ]} (Y_{t_{r}})$,\n\\end{itemize}\nthen $b_{k} (x)$ is given by the block\ndiagonal matrix\n\\[\nb_{k} (x) = \\left[ \\begin{array}{cccc}\nb_{t_{1}} (x) & & &\n\\\\\n& \\alpha^{-n (t_{1})} (b_{t_{2}}) (x) & &\n\\\\\n& &\n\\ddots &\n\\\\\n& & & \\alpha^{-[n (t_{1}) + \\cdots + n (t_{r - 1})]}\n(b_{t_{r}}) (x) \\end{array} \\right].\n\\]\n\\end{lem}\n\n\\begin{proof}\nThe proof is analogous (with appropriate changes to notation\nand exponents) to the proof of Lemma 11.3.18 of~\\cite{PhNotes}.\n\\end{proof}\n\nWe are now in position to give a decomposition of\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nas a recursive subhomogeneous algebra over $D$.\n\n\\begin{thm}\\label{AyRSHA}\nLet $X$ be a compact Hausdorff space, let $h \\colon X \\to X$ be a minimal homeomorphism,\nlet $D$ be a unital C*-algebra,\nand let $\\alpha \\in {\\mathrm{Aut}} (C (X, D))$ lie over~$h$.\nLet $Y \\subset X$ be closed with ${\\mathrm{int}} (Y) \\neq \\varnothing$.\nAdopt Notation \\ref{ModifiedRokhlinTower} and the notation of\nProposition \\ref{AyHom}.\nThen the homomorphism\n\\[\n\\gamma \\colon C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y\n \\to \\bigoplus_{k = 0}^{l} C (Y_{k}, M_{n (k)} (D))\n\\]\nof Proposition \\ref{AyHom}\ninduces an isomorphism of $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nwith the recursive subhomogeneous algebra over~$D$ defined,\nin the notation of Definition \\ref{RSHAMachinery}, as follows:\n\\begin{enumerate}\n\\item\\label{9223_AyRSHA_1}\n$l$ and $n (0),n (1),\\ldots,n (l)$ are as in Notation\n\\ref{ModifiedRokhlinTower}.\n\\item\\label{9223_AyRSHA_2}\n$X_{k} = Y_{k}$ for $0 \\leq k \\leq l$.\n\\item\\label{9223_AyRSHA_3}\n$X_{k}^{(0)} = Y_{k} \\cap \\bigcup_{j= 0}^{k - 1} Y_{j}$.\n\\item\\label{9223_AyRSHA_4}\nFor $x \\in X_{k}^{(0)}$ and $b = (b_{0}, b_{1},\\ldots, b_{k - 1})$\nin the image in $\\bigoplus_{j= 0}^{k - 1} C (Y_{j}, M_{n (j)} (D))$\nof the $k - 1$ stage algebra\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y^{(k - 1)}$, whenever\n\\[\nx \\in (Y_{k} \\setminus Y_{k}^{\\circ}) \\cap Y_{t_{1}} \\cap\nh^{-n (t_{1})} (Y_{t_{2}}) \\cap \\cdots \\cap h^{-[n (t_{1}) + \\cdots +\nn (t_{r - 1}]} (Y_{t_{r}})\n\\]\nwith\n$n (t_{1}) + n (t_{2}) + \\cdots + n (t_{r}) = n (k)$,\nthen $\\varphi_{k} (b (x))$ is given by the block\ndiagonal matrix\n\\[\n\\varphi_{k} (b (x))\n = \\left[ \\begin{array}{cccc}\nb_{t_{1}} (x) & & &\n\\\\\n& \\alpha^{-n (t_{1})} (b_{t_{2}}) (x) & &\n\\\\\n& & \\ddots &\n\\\\\n& & & \\alpha^{-[n (t_{1}) + \\cdots + n (t_{r - 1})]}\n(b_{t_{r}}) (x) \\end{array} \\right].\n\\]\n\\item\\label{9223_AyRSHA_5}\n$\\rho_{k}$ is the restriction map.\n\\end{enumerate}\nThe topological dimension of this decomposition is $\\dim (X)$, and\nthe standard representation of\n$\\sigma \\bigl( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y \\bigr)$\nis the inclusion map in\n$\\bigoplus_{k = 0}^{l} C (Y_{k}, M_{n (k)} (D))$.\n\\end{thm}\n\n\\begin{proof}\nThe proof is analogous to the proof of Theorem 11.3.19\nof~\\cite{PhNotes}, while the proof of the statement regarding the\ntopological dimension is the same as for the proof given\nin~\\cite{PhNotes} for the part of Theorem 11.3.14 that is not\nincluded in Theorem 11.3.19.\n\\end{proof}\n\n\\section{Structural properties of the orbit breaking subalgebra\n and the crossed product}\\label{StructAY}\n\nIn this section,\nwe give some general theorems on the structure of C*-algebra{s}\nof the form $C^* \\big( {\\mathbb{Z}}, \\, C_0 (X, D), \\, \\alpha \\big)$,\nwhich we derive from results on the structure of\norbit breaking subalgebras.\nWe give many explicit examples in Section~\\ref{Sec_InfRed}.\n\n\\begin{cnv}\\label{N_9215_AAY_ntn}\nIn this section,\nas in Definition~\\ref{D_4Y12_OrbSubalg}\nbut with additional restrictions,\n$X$ will be an infinite compact metric space,\n$h \\colon X \\to X$ will be a minimal homeomorphism,\n$D$ will be a simple unital C*-algebra,\nand $\\alpha \\in {\\mathrm{Aut}} ( C (X, D))$\nwill be an automorphism which lies over~$h$.\n\\end{cnv}\n\nFor a few results, simplicity is not needed.\n\n\\begin{prp}\\label{AyRSHADirLim}\nAdopt Convention~\\ref{N_9215_AAY_ntn}.\nFor any nonempty closed set $Y \\subset X$,\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$ is a direct limit of\nrecursive subhomogeneous algebras over~$D$\nof the form $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_{Y_n}$\nfor closed subsets $Y_n \\S X$ with ${\\mathrm{int}} (Y_n) \\neq \\varnothing$.\n\\end{prp}\n\n\\begin{proof}\nGiven $Y \\subset X$ closed, choose a sequence\n$(Y_{n})_{n = 1}^{\\infty}$ of closed subsets of $X$\nwith ${\\mathrm{int}} (Y_{n}) \\neq \\varnothing$ and\n$Y_{n + 1} \\subset Y_{n}$ for $n \\in {\\mathbb{Z}}_{> 0}$,\nand with $\\bigcap_{n = 1}^{\\infty} Y_{n} = Y$.\nThat $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$ is\na direct limit of separable recursive subhomogeneous algebras\nover $D$ follows by applying Theorem~\\ref{AyRSHA} and\nProposition~\\ref{AyDirLim}.\n\\end{proof}\n\n\\begin{prp}\\label{AYDStable}\nAdopt Convention~\\ref{N_9215_AAY_ntn}.\nLet $\\mathcal{E}$ be a separable unital C*-algebra{}\nthat is strongly selfabsorbing in the sense of \\cite{TomsWinter2}.\nAssume that $D$ is separable and $\\mathcal{E}$-stable.\nLet $Y \\subset X$ be closed and nonempty.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$ is $\\mathcal{E}$-stable.\n\\end{prp}\n\n\\begin{proof}\nIf ${\\mathrm{int}} (Y) \\neq \\varnothing$,\nthis follows immediately from Theorem \\ref{AyRSHA} and\nProposition \\ref{DStableRSHA}.\nThe case of a general nonempty closed set\nnow follows from Proposition~\\ref{AyRSHADirLim} above\nand Corollary 3.4 of~\\cite{TomsWinter2}.\n\\end{proof}\n\nIn particular, we get ${\\mathcal{Z}}$-stability.\n\n\\begin{cor}\\label{EStableGivesZStable}\nAdopt Convention~\\ref{N_9215_AAY_ntn}.\nLet $\\mathcal{E}$ be a separable unital C*-algebra{}\nthat is strongly selfabsorbing in the sense of \\cite{TomsWinter2}.\nAssume that $D$ is separable and $\\mathcal{E}$-stable.\nLet $Y \\subset X$ be closed and nonempty.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$ is ${\\mathcal{Z}}$-stable.\n\\end{cor}\n\n\\begin{proof}\nStrongly selfabsorbing C*-algebra{s} are ${\\mathcal{Z}}$-stable\nby Theorem~3.1 of~\\cite{WinterZStab},\nand $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$ is $\\mathcal{E}$-stable by\nProposition~\\ref{AYDStable}, so the conclusion is immediate.\n\\end{proof}\n\nWe now turn to ${\\mathcal{Z}}$-stability\nof the crossed product $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$.\nWhen $D$ is nuclear, it can be obtained from the results of\n\\cite{ABP-Zstab}.\nHowever, if nuclearity is not assumed, then the\nmain conclusion of that paper only implies\nthat $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is tracially\n${\\mathcal{Z}}$-stable (in the sense of~\\cite{HirOr}).\n\n\\begin{thm}\\label{ATraciallyZStable}\nAdopt Convention~\\ref{N_9215_AAY_ntn}.\nAssume that $D$ is a simple separable unital ${\\mathcal{Z}}$-stable\nC*-algebra{} which has a quasitrace.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nis tracially ${\\mathcal{Z}}$-stable.\nIf, in addition, $D$ is nuclear,\nthen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is ${\\mathcal{Z}}$-stable.\n\\end{thm}\n\n\\begin{proof}\nTheorem 3.3 of~\\cite{HirOr} implies that $D$ has strict comparison\nof positive elements, and so\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis a centrally large subalgebra\nof $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ by\nCorollary \\ref{AyCentLargeConds}(\\ref{9212_AyCentLargeConds_StrCmp}).\nAlso, $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis ${\\mathcal{Z}}$-stable by Corollary~\\ref{EStableGivesZStable},\nso Theorem~2.2 of~\\cite{ABP-Zstab} implies\nthat $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is\ntracially ${\\mathcal{Z}}$-stable.\nIf $D$ is nuclear, then it is ${\\mathcal{Z}}$-stable\nby Theorem 4.1 of~\\cite{HirOr},\nand $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is nuclear.\nSo ${\\mathcal{Z}}$-stability of $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nfollows from Theorem 2.3\nof~\\cite{ABP-Zstab}.\n\\end{proof}\n\n\\begin{cor}\nAdopt Convention~\\ref{N_9215_AAY_ntn}.\nLet $D$ be a simple separable unital nuclear ${\\mathcal{Z}}$-stable\nC*-algebra{} which has a quasitrace.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ has nuclear\ndimension at most 1.\n\\end{cor}\n\n\\begin{proof}\nThis follows immediately from Theorem~\\ref{ATraciallyZStable} above\nand Theorem~B of~\\cite{CETWW}.\n\\end{proof}\n\nCorollary~\\ref{EStableGivesZStable} and its consequences\nrequire no\nassumption about the underlying dynamical system other than minimality.\nIn particular, if $D$ is ${\\mathcal{Z}}$-stable then\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$ is\n${\\mathcal{Z}}$-stable when $X$ is infinite dimensional, and even when\n$\\mathrm{mdim} (X, h)$ (as in Notation~\\ref{N_0108_DimX}) is strictly positive.\nIf moreover $D$ is nuclear\nthen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is\n${\\mathcal{Z}}$-stable.\nThus, crossed products of the form\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ are\n${\\mathcal{Z}}$-stable\nwhenever $h$ is minimal and $D$ is\nsimple, separable, unital, nuclear, and ${\\mathcal{Z}}$-stable,\nregardless of anything else about the underlying dynamics.\n\nIf $D$ is not assumed to be ${\\mathcal{Z}}$-stable,\nthen we must make assumptions about the dynamical system\n$(X, h)$ besides just minimality to expect\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$ to have\ntractable structure.\nThe hypotheses in Proposition~\\ref{TsrRR0A_Y},\nTheorem~\\ref{P_9Z26_XisCantor},\nProposition~\\ref{P_9Z21_dim1},\nand Theorem~\\ref{T_9222_Tsr1AndRR0_CrPrd}\n(which also include conditions on~$D$)\nare surely much stronger than needed.\nThey have the advantage that the proofs can easily by obtained from\nresults already in the literature.\n\n\\begin{prp}\\label{TsrRR0A_Y}\nAdopt Convention~\\ref{N_9215_AAY_ntn},\nand assume that $X$ is the Cantor set.\nLet $Y \\subset X$ be closed and nonempty.\n\\begin{enumerate}\n\\item\\label{TsrRR0A_Y_tsr}\nIf ${\\mathrm{tsr}} (D) = 1$, then\n${\\mathrm{tsr}} \\bigl( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y \\bigr) = 1$.\n\\item\\label{TsrRR0A_Y_RR}\nIf ${\\mathrm{RR}} (D) = 0$, then\n${\\mathrm{RR}} \\bigl( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y \\bigr) = 0$.\n\\end{enumerate}\n\\end{prp}\n\n\\begin{proof}\nFirst suppose that ${\\mathrm{int}} (Y) \\neq \\varnothing$.\nThen, from Lemma~\\ref{CantorSetCase} and the remark immediately\nafter it,\nwe see there are are AF algebras $B_{0}, B_{1}, \\ldots, B_{k}$\nsuch that\n\\[\nC^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y\n \\cong \\bigoplus_{k = 0}^{l} B_{k} \\otimes D.\n\\]\nThe result follows immediately.\n\nSince stable rank one and real rank zero are preserved by direct limits,\nthe general case now follows from\nProposition~\\ref{AyRSHADirLim}.\n\\end{proof}\n\n\\begin{thm}\\label{P_9Z26_XisCantor}\nAdopt Convention~\\ref{N_9215_AAY_ntn},\nand assume that $X$ is the Cantor set.\nFurther assume that one of the following conditions holds:\n\\begin{enumerate}\n\\item\\label{9217_GetTsr1_AppInn}\nAll elements of $\\set{ \\alpha_{x} \\colon x \\in X} \\subset {\\mathrm{Aut}} (D)$\nare approximately inner.\n\\item\\label{9217_GetTsr1_StrCmp}\n$D$ has strict comparison of positive elements.\n\\item\\label{9217_GetTsr1_SP}\n$D$ has property (SP) and the order on projections over~$D$\nis determined by quasitraces.\n\\item\\label{9217_GetTsr1_CptGen}\nThe set\n$\\bigl\\{ \\alpha_x^n \\colon {\\mbox{$x \\in X$ and $n \\in {\\mathbb{Z}}$}} \\bigr\\}$\nis contained in a subset of ${\\mathrm{Aut}} (D)$\nwhich is compact in the topology\nof pointwise convergence in the norm on~$D$.\n\\end{enumerate}\nIf ${\\mathrm{tsr}} (D) = 1$, then\n${\\mathrm{tsr}} \\bigl( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr) \\bigr) = 1$,\nand if also ${\\mathrm{RR}} (D) = 0$, then\n${\\mathrm{RR}} \\bigl( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr) \\bigr) = 0$.\n\\end{thm}\n\nThe action in this theorem has the Rokhlin property.\nHowever, it seems to be unknown whether the crossed product of\neven a\nsimple unital C*-algebra{} with stable rank one by a Rokhlin automorphism\nstill has stable rank one.\nSee Problem~\\ref{Pb_9922_RokhlinTsr} and the discussion afterwards.\n\nWe also emphasize that the conclusion for real rank zero\nis only valid if stable rank one is also assumed.\nThe reason is the stable rank one hypothesis\nin Theorem~6.4 of~\\cite{ArPh}.\nThat theorem likely holds without stable rank one,\nbut this has not yet been proved.\nFor now, this is no great loss,\nsince, other than purely infinite simple C*-algebra{s}\n(which are covered by our Theorem~\\ref{T_0101_PICase}),\nthere are no known examples of simple unital C*-algebra{s} which have\nreal rank zero but not stable rank one.\n\n\\begin{proof}[Proof of Theorem~\\ref{P_9Z26_XisCantor}]\nChoose any one point subset $Y \\S X$.\nSince ${\\mathrm{tsr}} (D) = 1$,\nProposition \\ref{TsrRR0A_Y}(\\ref{TsrRR0A_Y_tsr}) implies\nthat the subalgebra $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\n(see Definition~\\ref{D_4Y12_OrbSubalg})\nhas stable rank one.\nIf also ${\\mathrm{RR}} (D) = 0$,\nthen Proposition \\ref{TsrRR0A_Y}(\\ref{TsrRR0A_Y_RR}) implies\n${\\mathrm{tsr}} \\big( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y \\big) = 1$.\nNow $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y$\nis a centrally large subalgebra of\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nby the appropriate part of Corollary \\ref{AyCentLargeConds}.\nSo ${\\mathrm{tsr}} \\big( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr) \\big) = 1$\nby Theorem 6.3 of~\\cite{ArPh},\nand if ${\\mathrm{RR}} (D) = 0$, then Theorem~6.4 of~\\cite{ArPh}\ngives ${\\mathrm{RR}} \\big( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr) \\big) = 0$.\n\\end{proof}\n\nWith more restrictive assumptions on~$D$,\nenough is known to get results when $\\dim (X) = 1$.\n\n\\begin{prp}\\label{P_9027_FibPrdT_sr1}\nLet $A$, $B$, and $C$ be C*-algebra{s},\nand let $\\varphi \\colon A \\to C$ and $\\psi \\colon B \\to C$ be homomorphism{s}.\nLet $D = A \\oplus_{C, \\varphi, \\psi} B$ (Definition~\\ref{pullback}).\nAssume that $\\psi$ is surjective,\nand that $A$ and $B$ have stable rank one.\nThen $D$ has stable rank one.\n\\end{prp}\n\n\\begin{proof}\nSet $J = {\\operatorname{Ker}} (\\psi)$.\nWe have a commutative diagram with exact rows (maps defined below):\n\\[\n\\begin{CD}\n0 @>>> J @>{\\lambda}>> D @>{\\pi}>> A @>>> 0 \\\\\n& & @VV{{\\mathrm{id}}_B}V @VV{\\sigma}V @VV{\\gamma}V \\\\\n0 @>>> J @>{\\iota}>> B @>{\\kappa}>> B\/J @>>> 0.\n\\end{CD}\n\\]\nHere $\\iota$ is the inclusion of $J$ in~$B$\nand $\\kappa$ is the quotient map,\n$\\pi \\colon D \\to A$ and $\\sigma \\colon D \\to B$\nare the restrictions of the coordinate projections\n$(a, b) \\to a$ and $(a, b) \\to b$,\nand $\\lambda (b) = (b, 0)$ for $b \\in J$.\nFor $a \\in A$,\ndefine $\\gamma (a)$ by\nfirst choosing $b \\in B$ such that $\\psi (b) = \\varphi (a)$,\nand then setting $\\gamma (a) = \\kappa (b)$.\nIt is easy to check that $\\gamma$ is well defined,\nand that the diagram commutes.\nAll parts of exactness are immediate except that\nsurjectivity of $\\pi$ follows from\nsurjectivity of~$\\psi$.\n\nThe algebra $A$ has stable rank one by hypothesis,\nand $J$ has stable rank one by Theorem~4.4 of~\\cite{Rffl6}.\nNext,\n${\\mathrm{ind}} \\colon K_1 (B \/ J) \\to K_0 (J)$ is the zero map\nby Corollary~2 of~\\cite{Nag},\nso $K_0 (\\iota)$ is injective by the long exact sequence in\nK-theory for the bottom row.\nSince $K_0 (\\iota) = K_0 (\\sigma) \\circ K_0 (\\lambda)$,\nit follows that $K_0 (\\lambda)$ is injective.\nTherefore ${\\mathrm{ind}} \\colon K_1 (A) \\to K_0 (J)$ is the zero map\nby the long exact sequence in\nK-theory for the top row.\nNow $D$ has stable rank one by Corollary~2 of~\\cite{Nag}.\n\\end{proof}\n\nWe state for convenient reference the result on the stable rank\nof direct limits.\n\n\\begin{lem}\\label{L_0111_tsr_lim}\nLet $(A_{\\lambda})_{\\lambda \\in \\Lambda}$ be a direct system of C*-algebra{s}.\nThen\n${\\mathrm{tsr}} \\bigl( \\varinjlim_{\\lambda} A_{\\lambda} \\bigr)\n \\leq \\liminf_{\\lambda} {\\mathrm{tsr}} (A_{\\lambda})$.\n\\end{lem}\n\nWe do not assume that the maps of the direct system are injective.\n\n\\begin{proof}[Proof of Lemma~\\ref{L_0111_tsr_lim}]\nWhen $\\Lambda = {\\mathbb{Z}}_{> 0}$,\nthis is Theorem~5.1 of~\\cite{Rffl6}.\nThe proof for general direct limits is the same.\n\\end{proof}\n\n\\begin{prp}\\label{P_9027_CXD_Tsr1}\nLet $D$ be a simple unital C*-algebra{}\nwith stable rank one and real rank zero,\nand suppose that $K_1 (D) = 0$.\nLet $X$ be a compact Hausdorff space{} with $\\dim (X) \\leq 1$.\nThen $C (X, D)$ has stable rank one.\n\\end{prp}\n\n\\begin{proof}\nIf $X$ is a point, this is trivial.\nIf $X = [0, 1]$,\nit follows from Theorem~4.3 of~\\cite{NagOsPh2001}.\nIf $X$ is a one dimensional finite complex,\nthe result follows from these facts\nand Proposition~\\ref{P_9027_FibPrdT_sr1}\nby induction on the number of cells.\n\nNow suppose $X$ is a compact metric space.\nBy Corollary 5.2.6 of~\\cite{SkiK},\nthere is an inverse system $(X_n)_{n \\in {\\mathbb{Z}}_{\\geq 0}}$\nof one dimensional finite complexes\nsuch that $X \\cong \\varprojlim_n X_n$.\nThen $C (X, D) \\cong \\varinjlim_n C (X_n, D)$.\nSince $C (X_n, D)$ has stable rank one for all $n \\in {\\mathbb{Z}}_{\\geq 0}$,\nit follows from Lemma~\\ref{L_0111_tsr_lim}\nthat $C (X, D)$ has stable rank one.\n\nFinally,\nlet $X$ be a general compact Hausdorff space{} with $\\dim (X) \\leq 1$.\nBy Theorem~1 of~\\cite{Mdsc}\n(in Section~3 of that paper),\nthere is an inverse system $(X_{\\lambda})_{\\lambda \\in \\Lambda}$\nof one dimensional compact metric space{s}\nsuch that $X \\cong \\varprojlim_{\\lambda} X_{\\lambda}$.\nIt now follows that $C (X, D)$ has stable rank one\nby the same reasoning as in the previous paragraph.\n\\end{proof}\n\n\\begin{prp}\\label{P_9027_Rsha_Tsr1}\nLet $D$ be a simple unital C*-algebra.\nLet $R$ be a recursive subhomogeneous C*-algebra{} over~$D$\n(Definition~\\ref{RSHA})\nwith topological dimension at most~$1$\n(Definition \\ref{RSHAMachinery}(\\ref{9223_RSHAMachinery_6})).\nIf $D$ has stable rank one and real rank zero,\nand $K_1 (D) = 0$,\nthen $R$ has stable rank one.\n\\end{prp}\n\n\\begin{proof}\nThe proof is the same as that of Proposition~\\ref{DStableRSHA},\nexcept using Proposition~\\ref{P_9027_FibPrdT_sr1}\nin place of Lemma~\\ref{DStablePullback},\nand using Proposition~\\ref{P_9027_CXD_Tsr1}.\n\\end{proof}\n\n\\begin{prp}\\label{P_9Z21_dim1}\nAdopt Convention~\\ref{N_9215_AAY_ntn},\nand assume that $\\dim (X) \\leq 1$,\nthat $D$ has stable rank one and real rank zero,\nand that $K_1 (D) = 0$.\nLet $Y \\subset X$ be closed and nonempty.\nThen\n${\\mathrm{tsr}} \\bigl( C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)_Y \\bigr) = 1$.\n\\end{prp}\n\n\\begin{proof}\nFirst suppose that ${\\mathrm{int}} (Y) \\neq \\varnothing$.\nThen the recursive subhomogeneous algebra over~$D$\nin Theorem~\\ref{AyRSHA}\nhas base spaces $Y_k$ which,\nfollowing Notation~\\ref{ModifiedRokhlinTower},\nare closed subsets of~$X$.\nBy Proposition 3.1.3 of~\\cite{Prs},\nthey therefore have covering dimension at most~$1$.\nThe result now follows from Proposition~\\ref{P_9027_Rsha_Tsr1}.\n\nStable rank one is preserved by direct limits\n(Lemma~\\ref{L_0111_tsr_lim}),\nso the general case now follows from\nProposition~\\ref{AyRSHADirLim}.\n\\end{proof}\n\nThe obvious examples of minimal homeomorphism{s}\nof one dimensional spaces are irrational rotations on~$S^1$,\nbut there are others.\nSee, for example, \\cite{GjdJhn},\nand the work on minimal homeomorphism{s} of the product of the Cantor set\nand the circle in \\cite{LnMti1} and~\\cite{LnMti3}.\n\n\\begin{thm}\\label{T_9222_Tsr1AndRR0_CrPrd}\nAdopt Convention~\\ref{N_9215_AAY_ntn},\nand assume that $\\dim (X) \\leq 1$,\nthat $D$ has stable rank one and real rank zero,\nand that $K_1 (D) = 0$.\nFurther assume that one of the following conditions holds:\n\\begin{enumerate}\n\\item\\label{9927_AyCentLargeConds_AppInn}\nAll elements of $\\set{ \\alpha_{x} \\colon x \\in X} \\subset {\\mathrm{Aut}} (D)$\nare approximately inner.\n\\item\\label{9927_AyCentLargeConds_SP}\nThe order on projections over~$D$\nis determined by quasitraces.\n\\item\\label{9927_AyCentLargeConds_CptGen}\nThe set\n$\\bigl\\{ \\alpha_x^n \\colon {\\mbox{$x \\in X$ and $n \\in {\\mathbb{Z}}$}} \\bigr\\}$\nis contained in a subset of ${\\mathrm{Aut}} (D)$\nwhich is compact in the topology\nof pointwise convergence in the norm on~$D$.\n\\end{enumerate}\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nhas stable rank one.\n\\end{thm}\n\n\\begin{proof}\nCombine Theorem 6.3 of~\\cite{ArPh},\nCorollary~\\ref{AyCentLargeConds},\nand Proposition~\\ref{P_9027_Rsha_Tsr1}.\n\\end{proof}\n\n\\section{Minimality of products of Denjoy homeomorphisms\n of the Cantor set}\\label{Sec_9929_MinProd}\n\n\\indent\nFor use in examples in the next section,\nwe prove minimality of products of Denjoy homeomorphisms\nwith rationally independent rotation numbers.\nWe have not found this result in the literature.\nThe method we use here is not the traditional approach\nto this kind of problem,\nbut we hope it will be useful elsewhere.\n\n\\begin{ntn}\\label{N_9X04_Rotation}\nFor $\\theta \\in {\\mathbb{R}}$ we let $r_{\\theta} \\colon S^1 \\to S^1$\ndenote the rotation by $2 \\pi \\theta$,\nthat is, $r_{\\theta} (\\zeta) = \\exp (2 \\pi i \\theta) \\zeta$\nfor $\\zeta \\in S^1$.\n\\end{ntn}\n\n\n\\begin{prp}\\label{P_4X04_ProdRotation}\nLet $I$ be a set,\nand let $(\\theta_i)_{i \\in I}$ be a family of elements of~${\\mathbb{R}}$.\nSuppose that $1$ and the numbers $\\theta_i$ for $i \\in I$ are\nlinearly independent over~${\\mathbb{Q}}$.\nThen the product $r \\colon (S^1)^{I} \\to (S^1)^{I}$\nof the rotations $r_{\\theta_i}$ is minimal.\n\\end{prp}\n\n\\begin{proof}\nIf $I$ is finite,\nthis is Proposition 1.4.1 of~\\cite{KtHsb}.\nThe general case follows from the finite case by taking the inverse\nlimit of the products over finite subsets of~$I$,\nby the discussion after Proposition II.4 of~\\cite{Frtg}.\nThe proof is is written for actions of the semigroup~${\\mathbb{Z}}_{\\geq 0}$\nand for countable inverse limits,\nbut the proof for actions of ${\\mathbb{Z}}$ and arbitrary inverse limits\nis the same.\n\\end{proof}\n\n\\begin{dfn}\\label{D_9929_StrOpen}\nLet $X$ and $Y$ be topological spaces,\nand let $g \\colon Y \\to X$.\nLet $y \\in Y$.\nWe say that $g$ is {\\emph{strictly open}} at~$y$\nif\n\\[\n\\bigl\\{ g^{-1} (U) \\colon\n {\\mbox{$U \\S X$ is open and $g (y) \\in U$}} \\bigr\\}\n\\]\nis a neighborhood base for $y$ in~$Y$.\n\\end{dfn}\n\nSome of the theory below works without assuming $g$ is continuous,\nso we do not require that neighborhood{s} be open.\n\nThe next lemma gives an alternate interpretation of strict openness.\n\n\\begin{lem}\\label{L_9929_WhenStrOpen}\nIn the situation in Definition~\\ref{D_9929_StrOpen},\nassume that $g$ is continuous{} and surjective, $Y$ is compact,\nand $X$ and $Y$ are Hausdorff.\nThen $g$ is strictly open at~$y$ if and only if{}\nthe following conditions hold:\n\\begin{enumerate}\n\\item\\label{Item_D_9929_Inj}\n$g^{-1} (\\{ g (y) \\}) = \\{ y \\}$.\n\\item\\label{Item_D_9929_Open}\nFor every open set $V \\S Y$ with $y \\in V$,\nthe set $g (V)$ is a neighborhood{} of $g (y)$.\n\\end{enumerate}\n\\end{lem}\n\nThe lemma is false without compactness of~$Y$.\nTake\n\\[\nX = [0, 1]\n\\qquad {\\mbox{and}} \\qquad\nY = \\{ (0, 0) \\} \\cup \\bigl( (0, 1]\\times [0, 1] \\bigr) \\S [0, 1]^2,\n\\]\nlet $g$ be projection to the first coordinate,\nand take $y = (0, 0)$.\nThen (\\ref{Item_D_9929_Inj}) and~(\\ref{Item_D_9929_Open}) hold,\nbut $g$ is not strictly open at~$y$.\n\n\\begin{proof}[Proof of Lemma~\\ref{L_9929_WhenStrOpen}]\nAssume that $g$ is strictly open at~$y$.\nTo prove~(\\ref{Item_D_9929_Inj}),\nlet $z \\in Y \\setminus \\{ y \\}$,\nchoose $V \\S Y$ open with $y \\in V$ and $z \\not\\in V$,\nand choose $U \\S X$ open such that $g (y) \\in U$\nand $g^{-1} (U) \\S V$.\nThen $z \\not\\in g^{-1} (U)$, so $g (z) \\neq g (y)$.\n\nTo prove~(\\ref{Item_D_9929_Open}),\nlet $V \\S Y$ be open with $y \\in V$.\nChoose $U \\S X$ open such that $g (y) \\in U$\nand $g^{-1} (U) \\S V$.\nSurjectivity of~$g$ implies $U \\S g ( g^{-1} (U) ) \\S g (V)$,\nso $g (V)$ is a neighborhood{} of $g (y)$.\n\nFor the converse,\nassume~(\\ref{Item_D_9929_Open}),\nand suppose that $g$ is not strictly open at~$y$.\nLet ${\\mathcal{N}}$ be the set of open subsets of~$X$ which contain $g (y)$,\nordered by reverse inclusion.\nSince $g$ is continuous,\nfailure of strict openness means we can\nchoose $V \\S Y$ open with $y \\in V$\nand such that for every $U \\in {\\mathcal{N}}$,\nwe have $g^{-1} (U) \\not\\S V$.\nSo there is $z_U \\in Y$ such that $z_U \\not\\in V$ but $g (z_U) \\in U$.\nThe net $(z_U)_{U \\in {\\mathcal{N}}}$\nsatisfies $g (z_U) \\to g (y)$.\nChoose a subnet $(t_{\\lambda})_{\\lambda \\in \\Lambda}$\nsuch that $t = \\lim_{\\lambda} t_{\\lambda}$ exists.\nThen $t \\not\\in V$.\nAlso $g (t_{\\lambda}) \\to g (t)$ by continuity and $g (t_{\\lambda}) \\to g (y)$\nby construction, so $g (t) = g (y)$.\nWe have contradicted~(\\ref{Item_D_9929_Inj}).\n\\end{proof}\n\n\\begin{lem}\\label{L_9929_ExtMin}\nLet $X$ and $Y$ be compact Hausdorff space{s},\nand let $h \\colon X \\to X$ and $k \\colon Y \\to Y$ be homeomorphism{s}.\nLet $g \\colon Y \\to X$ be continuous and surjective,\nand satisfy $g \\circ k = h \\circ g$.\nSuppose that $h$ is minimal and there is a dense subset $S \\S Y$\nsuch that $g$ is strictly open at every point of~$S$.\nThen $k$ is minimal.\n\\end{lem}\n\n\\begin{proof}\nIt is enough to show that for every $y \\in Y$ and every\nnonempty open set $V \\S Y$,\nthere is $n \\in {\\mathbb{Z}}$ such that $k^n (y) \\in V$.\nChoose $z \\in S \\cap V$.\nSince $g$ is strictly open at~$z$,\nthere is an open set $U \\S X$ such that $g (z) \\in U$\nand $g^{-1} (U) \\S V$.\nSince $h$ is minimal,\nthere is $n \\in {\\mathbb{Z}}$ such that $h^n (g (y)) \\in U$.\nThen $g (k^n (y)) \\in U$,\nso $k^n (y) \\in g^{-1} (U) \\S V$.\n\\end{proof}\n\nIt is well known that Lemma~\\ref{L_9929_ExtMin} fails\nwithout strict openness,\neven if one assumes that $g^{-1} (x)$ is finite for all $x \\in X$\nand has only one point for a dense subset of points $x \\in X$.\nThe following example was suggested by B.~Weiss.\n\n\\begin{exa}\\label{E_9X04_NeedStrict}\nFix $\\theta \\in {\\mathbb{R}} \\setminus {\\mathbb{Q}}$.\nTake $X = S^1$ and take $h = r_{\\theta}$\n(rotation by $2 \\pi \\theta$; Notation~\\ref{N_9X04_Rotation}).\nDefine $Y \\S S^1 \\times [0, 1]$ to be\n\\[\nY = \\bigl( S^1 \\times \\{ 0 \\} \\bigr)\n \\cup \\left\\{ \\left( e^{2 \\pi i n \\theta}, \\, \\frac{1}{| n | + 1} \\right)\n \\colon n \\in {\\mathbb{Z}} \\right\\},\n\\]\nand define $k \\colon Y \\to Y$ to be\n$k (\\zeta, 0) = (r_{\\theta} (\\zeta), \\, 0)$ for $\\zeta \\in S^1$ and\n\\[\nk \\left( e^{2 \\pi i n \\theta}, \\, \\frac{1}{| n | + 1} \\right)\n = \\left( e^{2 \\pi i (n + 1) \\theta}, \\, \\frac{1}{| n + 1 | + 1} \\right)\n\\]\nfor $n \\in {\\mathbb{Z}}$.\nLet $g \\colon Y \\to X$ be projection to the first coordinate.\nThen $g \\circ k = h \\circ g$,\n$g$ is surjective,\n$h$ is minimal,\n$k$ is not minimal,\n$g^{-1} (x)$ has at most two points for every $x \\in X$,\nand $g^{-1} (x)$ has one point for all but countably many $x \\in X$.\n\\end{exa}\n\nThe following two easy lemmas will be used in the construction of\nexamples.\n\n\\begin{lem}\\label{L_9X04_Prod}\nLet $I$ be a nonempty set,\nand for $i \\in I$ let $g_i \\colon X_i \\to Y_i$\nand $y_i \\in Y_i$ be as in Definition~\\ref{D_9929_StrOpen},\nwith $g_i$ strictly open at~$y_i$.\nSet\n\\[\nX = \\prod_{i \\in I} X_i\n\\qquad {\\mbox{and}} \\qquad\nY = \\prod_{i \\in I} Y_i,\n\\]\nlet $g \\colon Y \\to X$ be given by $g (z) = ( g_i (z_i))_{i \\in I}$\nfor $z = (z_i)_{i \\in I} \\in Y$,\nand set $y = (y_i)_{i \\in I}$.\nThen $g$ is strictly open at~$y$.\n\\end{lem}\n\n\\begin{proof}\nIt is obvious that $g^{-1} ( \\{ g (y) \\} ) = \\{ y \\}$.\nNow let $V \\S Y$ be open with $y \\in V$.\nWe need an open set $U \\S X$\nsuch that $g (y) \\in U$, $g^{-1} (U)$ is a neighborhood{} of~$y$,\nand $g^{-1} (U) \\S V$.\nWithout loss of generality{} there are open subsets $V_i \\S Y_i$\nand a finite set $F \\S I$\nsuch that $V = \\prod_{i \\in I} V_i$,\n$V_i = Y_i$ for all $i \\in I \\setminus F$,\nand $y_i \\in V_i$ for all $i \\in I$.\nFor $i \\in F$ choose an open set $U_i \\S X_i$ with $g_i (y_i) \\in U_i$\nand such that $g_i^{-1} (U_i) \\S V_i$.\nSet $U_i = X_i$ for $i \\in I \\setminus F$.\nSince $g_i^{-1} (U_i)$ is a neighborhood{} of~$y_i$\nfor all $i \\in I$\nand $U_i = X_i$ for $i \\in I \\setminus F$,\nthe set $g^{- 1} (U) = \\prod_{i \\in I} g_i^{-1} (U_i)$\nis a neighborhood{} of~$y$.\nSo the set $U = \\prod_{i \\in I} U_i$ is the required set.\n\\end{proof}\n\n\\begin{lem}\\label{L_9X04_Subsp}\nIn the situation in Definition~\\ref{D_9929_StrOpen},\nassume that $g$ is strictly open at~$y$.\nLet $Y_0 \\S Y$ be a subspace such that $y \\in Y_0$.\nThen $g |_{Y_0}$ is strictly open at~$y$.\n\\end{lem}\n\n\\begin{proof}\nThe result is immediate from the definition of the subspace\ntopology.\n\\end{proof}\n\n\\begin{lem}\\label{L_9929_S1toS1}\nLet $g \\colon S^1 \\to S^1$ be continuous{} and surjective.\nSuppose that there are disjoint closed arcs\n$I_1, I_2, \\ldots \\S S^1$\nsuch that:\n\\begin{enumerate}\n\\item\\label{Item_L_9929_S1toS1_Const}\nFor $n \\in {\\mathbb{Z}}_{> 0}$, the function $g$ is constant on~$I_n$.\n\\item\\label{Item_L_9929_S1toS1_NI}\n$\\bigcup_{n = 1}^{\\infty} I_n$ is dense in~$S^1$.\n\\item\\label{Item_L_9929_S1toS1_Inj}\nWith $R = S^1 \\setminus \\bigcup_{n = 1}^{\\infty} I_n$,\nthe restriction $g |_R$ is injective.\n\\end{enumerate}\nThen $g$ is strictly open at every point of~$R$.\n\\end{lem}\n\n\\begin{proof}\nFor distinct $\\lambda, \\zeta \\in S^1$,\nwe denote the open, closed, and half open arcs from $\\lambda$ to~$\\zeta$\nusing interval notation:\n$(\\lambda, \\zeta)$, $[\\lambda, \\zeta]$, $[\\lambda, \\zeta)$, and $(\\lambda, \\zeta]$.\nFurther write $I_n = [\\beta_n, \\gamma_n]$ with $\\beta_n, \\gamma_n \\in S^1$,\nand set $B = \\{ \\beta_n \\colon n \\in {\\mathbb{Z}}_{> 0} \\}$\nand $C = \\{ \\gamma_n \\colon n \\in {\\mathbb{Z}}_{> 0} \\}$.\n\nThe set $g (S^1 \\setminus R)$ is countable\nand $g$ is surjective, so $S^1 \\setminus g (R)$ is countable.\n\nFirst,\nwe claim that if $E, F \\S S^1$ are disjoint,\nthen ${\\mathrm{int}} (g (E)) \\cap {\\mathrm{int}} (g (F)) = \\varnothing$.\nIf the claim is false,\nthen $g (E) \\cap g (F)$ is uncountable.\nNow $g (E) \\setminus g (E \\cap R) \\S g (S^1 \\setminus R)$\nwhich is countable, so $g (E) \\setminus g (E \\cap R)$ is countable.\nSimilarly $g (F) \\setminus g (F \\cap R)$ is countable.\nSo $g (E \\cap R) \\cap g (F \\cap R)$ is uncountable.\nThis contradicts $E \\cap F = \\varnothing$ and injectivity of $g |_R$.\n\nSecond, we claim that if $J \\S S^1$ is any nonempty open arc such that\nthere is no $n \\in {\\mathbb{Z}}_{> 0}$ with ${\\overline{J}} \\S I_n$,\nthen $R \\cap J \\neq \\varnothing$.\nSuppose the claim is false.\nWrite $J = (\\lambda, \\zeta)$ with $\\lambda, \\zeta \\in S^1$.\nDefine\n\\[\nS = \\begin{cases}\n \\varnothing & \\hspace*{1em} \\lambda \\not\\in R\n \\\\\n \\{ \\lambda \\} & \\hspace*{1em} \\lambda \\in R\n\\end{cases}\n\\qquad {\\mbox{and}} \\qquad\nT = \\begin{cases}\n \\varnothing & \\hspace*{1em} \\zeta \\not\\in R\n \\\\\n \\{ \\zeta \\} & \\hspace*{1em} \\zeta \\in R.\n\\end{cases}\n\\]\nThen ${\\overline{J}}$ is the disjoint union of the closed sets\n\\[\nS, \\, T, \\, {\\overline{J}} \\cap I_1, \\, {\\overline{J}} \\cap I_2, \\ldots.\n\\]\nSuppose there is $n \\in {\\mathbb{Z}}_{> 0}$ such that ${\\overline{J}} \\cap I_m = \\varnothing$\nfor all $m \\neq n$.\nSince ${\\overline{J}} \\not\\subset I_n$,\nthere is a nonempty open arc $L \\S J$\nsuch that $L \\cap I_n = \\varnothing$,\nwhence $L \\cap I_m = \\varnothing$ for all $m \\in {\\mathbb{Z}}_{> 0}$.\nThis contradicts~(\\ref{Item_L_9929_S1toS1_NI}).\nThus, at least two of\nthe sets ${\\overline{J}} \\cap I_n$ are nonempty.\nBut, according to\nTheorem~6 in Part III of Section~47 (in Chapter~5) of~\\cite{Krtw},\na connected compact Hausdorff space{} cannot be the disjoint union of closed subsets\n$E_1, E_2, \\ldots$ with at least two of them nonempty.\nThis contradiction proves the claim.\n\nThird, we claim that, under the same hypotheses, $R \\cap J$\nis infinite.\nAgain write $J = (\\lambda, \\zeta)$ with $\\lambda, \\zeta \\in S^1$.\nApply the previous claim\nto choose $\\rho_1 \\in R \\cap J$ and set $J_1 = (\\rho_1, \\zeta)$.\nThen $\\rho_1 \\not\\in \\bigcup_{n = 1}^{\\infty} I_n$\nso there is no $n \\in {\\mathbb{Z}}_{> 0}$ with ${\\overline{J_1}} \\S I_n$.\nApply the previous claim again, choosing $\\rho_2 \\in R \\cap J_1$,\nand showing that $J_2 = (\\rho_2, \\zeta)$\nsatisfies ${\\overline{J_1}} \\not\\S I_n$ for all $n \\in {\\mathbb{Z}}_{> 0}$.\nProceed inductively.\n\nFourth, we claim that if $\\lambda \\in R \\cup C$ and $\\zeta \\in R \\cup B$\nare distinct,\nthen $g \\bigl( (\\lambda, \\zeta) \\bigr)$ is an open arc.\nFor the proof, set $J = (\\lambda, \\zeta)$.\nWe can assume that $g (J) \\neq S^1$.\nSince $g (J)$ is connected, there are $\\mu, \\nu \\in S^1$\nsuch that $g (J)$ is one of\n$(\\mu, \\nu)$, $[\\mu, \\nu]$\n(in this case we allow $\\mu = \\nu$), $[\\mu, \\nu)$, or $(\\mu, \\nu]$.\nSuppose $\\mu \\in g (J)$; we will obtain a contradiction.\nChoose $\\rho \\in J$ such that $g (\\rho) = \\mu$.\nThere are two cases: $\\rho \\in R$ and $\\rho \\in \\bigcup_{n = 1}^{\\infty} I_n$.\n\nIn the first case, let $L_0$ and $L_1$ be the nonempty arcs\n$L_0 = (\\lambda, \\rho)$ and $L_1 = (\\rho, \\zeta)$.\nFor $j = 0, 1$,\nsince $\\rho \\in {\\overline{L_j}}$,\nthe arc $L_j$ satisfies the hypotheses of the second claim.\nSo $R \\cap L_j$ is infinite by the third claim.\nSince $g |_R$ is injective, $g (L_j)$ is infinite.\nClearly $\\mu \\in {\\overline{g (L_j)}} \\S [\\mu, \\nu]$.\nAlso, $g (L_j)$ is an arc by connectedness.\nSo there is $\\omega_j \\in (\\mu, \\nu)$ such that $(\\mu, \\omega_j) \\S g (L_j)$.\nThis implies ${\\mathrm{int}} ( g (L_0) ) \\cap {\\mathrm{int}} ( g (L_1) ) \\neq \\varnothing$,\ncontradicting the first claim.\n\nIn the second case,\nlet $n \\in {\\mathbb{Z}}_{> 0}$ be the integer such that $\\rho \\in I_n$.\nBoth $\\lambda \\in [\\beta_n, \\gamma_n)$ and $\\zeta \\in (\\beta_n, \\gamma_n]$\nare impossible, but $[\\beta_n, \\gamma_n] \\cap (\\lambda, \\zeta) \\neq \\varnothing$,\nso $[\\beta_n, \\gamma_n] \\subset (\\lambda, \\zeta)$.\nTherefore the arcs $L_0 = (\\lambda, \\beta_n)$ and $L_1 = (\\gamma_n, \\zeta)$\nare nonempty and do not intersect~$I_n$.\nIf there is $m \\in {\\mathbb{Z}}_{> 0}$ such that ${\\overline{L_0}} \\S I_m$,\nthen clearly $m \\neq n$,\nbut then $\\beta_n \\in I_m \\cap I_n$,\ncontradicting $I_m \\cap I_n = \\varnothing$.\nSo no such $m$ exists,\nwhence $L_0$ satisfies the hypotheses of the second claim.\nSimilarly $L_1$ satisfies the hypotheses of the second claim.\nSince $g$ is constant on~$I_n$,\nwe have $g (\\beta_n) = g (\\gamma_n) = g (\\rho) = \\mu$.\nNow we get a contradiction in the same way as in the first case.\n\nSo $\\mu \\not\\in g (J)$.\nSimilarly $\\nu \\not\\in g (J)$.\nThe claim is proved.\n\nFifth,\nwe claim that if $L \\S S^1$ is an open arc,\nthen $g^{-1} (L)$ is an open arc.\nTo prove the claim,\nsince $g^{-1} (L)$ is open,\nthere are an index set $S$\nand disjoint nonempty open arcs $J_s$ for $s \\in S$\nsuch that $g^{-1} (L) = \\bigcup_{s \\in S} J_s$,\nand there are $\\lambda_s, \\zeta_s \\in S^1$\nsuch that $J_s = (\\lambda_s, \\zeta_s)$.\n\nLet $s \\in S$.\nSuppose $\\lambda_s \\in I_n$ for some $n \\in {\\mathbb{Z}}_{> 0}$.\nSince $g$ is constant on~$I_n$,\neither $I_n \\cap g^{-1} (L) = \\varnothing$\nor $I_n \\S g^{-1} (L)$.\nIn the second case, $I_n \\cup (\\lambda_s, \\zeta_s)$\nis a connected subset of $g^{-1} (L)$.\nSince $(\\lambda_s, \\zeta_s)$ is a maximal connected subset of $g^{-1} (L)$,\nwe have $I_n \\S (\\lambda_s, \\zeta_s)$.\nThis contradicts $\\lambda_s \\in I_n$.\nSo the first case must apply, and therefore $\\lambda_s = \\gamma_n$.\nCombining this with the possibility\n$\\lambda_s \\not\\in \\bigcup_{n = 1}^{\\infty} I_n$,\nwe conclude that $\\lambda_s \\in R \\cup C$.\nSimilarly, $\\zeta_s \\in R \\cup B$.\nBy the previous claim,\n$g \\bigl( (\\lambda_s, \\zeta_s) \\bigr)$ is open.\n\nThe first claim now implies that the sets\n$g \\bigl( (\\lambda_s, \\zeta_s) \\bigr)$ are all disjoint.\nSince $L$ is connected,\nit follows that ${\\mathrm{card}} (S) = 1$,\nwhich implies the claim.\n\nTo prove the lemma,\nlet $\\eta \\in R$ and let $V \\S S^1$ be open with $\\eta \\in V$.\nWe need an open set $U \\S S^1$\nsuch that $g (\\eta) \\in U$ and $g^{-1} (U) \\S V$.\nChoose $\\lambda_0, \\zeta_0 \\in S^1$\nsuch that $\\eta \\in (\\lambda_0, \\zeta_0) \\S V$.\nUse the second claim to choose $\\lambda_1 \\in (\\lambda_0, \\eta) \\cap R$\nand $\\zeta_1 \\in (\\eta, \\zeta_0) \\cap R$.\nThen $U = g \\bigl( (\\lambda_1, \\zeta_1) \\bigr)$ is open by\nthe fourth claim,\nand by the fifth claim there are $\\lambda_2, \\zeta_2 \\in S^1$\nsuch that $g^{-1} (U) = (\\lambda_2, \\zeta_2)$.\nIf $\\lambda_1 \\in g^{-1} (U)$,\nthen $(\\lambda_2, \\lambda_1)$ is a nonempty open arc contained in $g^{-1} (U)$\nwith $\\lambda_1 \\in R$.\nThe third claim and injectivity of $g |_R$\nimply that $g \\bigl( (\\lambda_2, \\lambda_1) \\bigr)$\nis not a point.\nSince $g \\bigl( (\\lambda_2, \\lambda_1) \\bigr)$ is connected,\nit has nonempty interior.\nIt is contained in $U = g \\bigl( (\\lambda_1, \\zeta_1) \\bigr)$,\ncontradicting the first claim.\nSo $\\lambda_1 \\not\\in (\\lambda_2, \\zeta_2)$.\nSimilarly $\\zeta_1 \\not\\in (\\lambda_2, \\zeta_2)$.\nSince $\\eta \\in (\\lambda_2, \\zeta_2)$,\n$\\eta \\in (\\lambda_1, \\zeta_1)$, and $(\\lambda_2, \\zeta_2)$ is connected,\nit now follows that $(\\lambda_2, \\zeta_2) \\S (\\lambda_1, \\zeta_1)$.\n(In fact, we have equality, since $g \\bigl( (\\lambda_1, \\zeta_1) \\bigr) = U$.)\nThus\n\\[\ng^{-1} (U)\n = (\\lambda_2, \\zeta_2)\n \\S (\\lambda_1, \\zeta_1)\n \\S (\\lambda_0, \\zeta_0)\n \\S V.\n\\]\nThis completes the proof.\n\\end{proof}\n\nWe have not seen a term in the literature for the following\nclass of homeomorphism{s}.\n\n\\begin{dfn}\\label{D_9X04_RestrDj}\nA {\\emph{restricted Denjoy homeomorphism}}\nis a homeomorphism of the Cantor set which is conjugate to\nthe restriction and corestriction of a Denjoy homeomorphism of~$S^1$,\nin the sense of Definition~3.3 of~\\cite{PtSmdSk},\nto the unique minimal set of Proposition~3.4 of~\\cite{PtSmdSk}.\nThe {\\emph{rotation number}}\nof a restricted Denjoy homeomorphism is the rotation number,\nas at the beginning of Section~3 of~\\cite{PtSmdSk},\nof the Denjoy homeomorphism of~$S^1$\nof which it is the restriction;\nthis number is well defined by Remark~3 at the end of Section~3\nof~\\cite{PtSmdSk}.\n\\end{dfn}\n\n\\begin{lem}\\label{L_9X04_DjSemiCj}\nLet $X$ be the Cantor set,\nand let $k \\colon X \\to X$ be a restricted Denjoy homeomorphism\nwith rotation number~$\\theta$,\nas in Definition~\\ref{D_9X04_RestrDj}.\nThen there are a continuous{} surjective map $g \\colon X \\to S^1$\nand a dense subset $R \\S X$ such that\n$g$ is strictly open at~$x$ for every $x \\in R$\nand, following Notation~\\ref{N_9X04_Rotation},\n$g \\circ h = r_{\\theta} \\circ g$.\n\\end{lem}\n\n\\begin{proof}\nWe may assume that $k_0 \\colon S^1 \\to S^1$ is a Denjoy homeomorphism\nas in Definition~3.3 of~\\cite{PtSmdSk},\nthat $X \\S S^1$ is the unique minimal set for~$k_0$,\nand that $k$ is the restriction and corestriction of $k_0$ to~$X$.\nCorollary~3.2 of~\\cite{PtSmdSk}\ngives $h \\colon S^1 \\to S^1$\nsuch that $h \\circ k_0 = r_{\\theta} \\circ h$.\nSet $g = h |_X$.\nSince $X \\S S^1$ is compact and invariant and $r_{\\theta}$ is minimal,\n$g$ must be surjective.\nThe discussion after Proposition~3.4 of~\\cite{PtSmdSk}\nimplies that $h$ satisfies the hypotheses of Lemma~\\ref{L_9929_S1toS1},\nwith $X = S^1 \\setminus \\bigcup_{n = 1}^{\\infty} {\\mathrm{int}} (I_n)$.\nMoreover, the set $R$ of Lemma~\\ref{L_9929_S1toS1}\nis nonempty by\nTheorem~6 in Part III of Section~47 (in Chapter~5) of~\\cite{Krtw},\nand $k$-invariant,\nhence dense in~$X$.\nBy Lemma~\\ref{L_9929_S1toS1} and Lemma~\\ref{L_9X04_Subsp},\n$g$ is strictly open at every point of~$R$.\n\\end{proof}\n\n\\begin{prp}\\label{P_4X04_ProdRstDj}\nLet $I$ be a set,\nlet $(\\theta_i)_{i \\in I}$ be a family of elements of~${\\mathbb{R}}$,\nand for $i \\in I$ let $k_i \\colon X_i \\to X_i$\nbe a restricted Denjoy homeomorphism with rotation number~$\\theta_i$,\nas in Definition~\\ref{D_9X04_RestrDj}.\nSuppose that $1$ and the numbers $\\theta_i$ for $i \\in I$ are\nlinearly independent over~${\\mathbb{Q}}$.\nThen the product $k \\colon \\prod_{i \\in I} X_i \\to \\prod_{i \\in I} X_i$\nof the homeomorphism{s} $k_i$ is minimal.\n\\end{prp}\n\n\\begin{proof}\nFor $i \\in I$,\nlet $g_i \\colon X_i \\to S^1$ be as in Lemma~\\ref{L_9X04_DjSemiCj}.\nBy Lemma~\\ref{L_9X04_DjSemiCj},\nthere is a dense subset $R_i \\S X_i$\nsuch that $g_i$ is strictly open at every point of~$R_i$.\nLet $g \\colon \\prod_{i \\in I} X_i \\to (S^1)^{I}$\nbe the product of the maps~$g_i$.\nLemma~\\ref{L_9X04_Prod} implies that\n$g$ is strictly open at every point of $\\prod_{i \\in I} R_i$,\nthis set is dense,\nand the product $r \\colon (S^1)^{I} \\to (S^1)^{I}$\nof the rotations $r_{\\theta_i}$ is minimal\nby Proposition~\\ref{P_4X04_ProdRotation}.\nSo $k$ is minimal\nby Lemma~\\ref{L_9929_ExtMin}.\n\\end{proof}\n\n\n\\section{Examples}\\label{Sec_InfRed}\n\n\\indent\nTo take full advantage of Proposition~\\ref{AyRSHADirLim},\nwe need information on the structure of simple\ndirect limits of recursive subhomogeneous algebras\nover the algebra $D$ which appears there.\nFor example,\nwe need conditions for such direct limits to have stable rank~$1$\n(as originally done for $D = {\\mathbb{C}}$ in~\\cite{DNNP})\nand to have real rank zero\n(as originally done for $D = {\\mathbb{C}}$ in~\\cite{BDR}).\nWe intend to study this problem in a future paper.\nEven without generalizations of those theorems,\nthree cases are accessible now.\nThey are $h \\colon X \\to X$ wild\n(for example, with strictly positive mean dimension,\nas in Notation~\\ref{N_0108_DimX})\nbut with $D$ being ${\\mathcal{Z}}$-stable;\n$X$ is the Cantor set\nand $D$ has stable rank one but is otherwise wild;\nand $X = S^1$,\n$h$ is an irrational rotation,\nand $D$ has stable rank one, real rank zero,\nand trivial $K_1$-group,\nbut is otherwise wild.\n\nWe also give examples for Theorem~\\ref{T_0101_PICase},\nalthough only for $G = {\\mathbb{Z}}$.\n\nIn the first type of example,\nif $X$ is finite dimensional{} then the action has a\nhigher dimensional Rokhlin property,\nand results on crossed products by such actions can be applied.\nAs far as we know, however,\nthere are no previously known general theorems\nwhich apply when $X$ is infinite dimensional.\nIn the second type of example,\nthe action has the Rokhlin property.\nSince the algebra is not simple,\nresults of~\\cite{OP1} can't be applied,\nand, when $D$ is not ${\\mathcal{Z}}$-stable\nand does not have finite nuclear dimension,\nTheorems 4.1 and~5.8 of~\\cite{HWZ} can't be applied.\nWe know of no previous general theorems which apply in this case.\nActions in the third type of example\nare further away from known results:\nthe situation is like for the second type of example,\nbut the Rokhlin property must be weakened to\nfinite Rokhlin dimension with commuting towers.\n\nSince we will use quasifree automorphisms of reduced C*-algebras\nof free groups\nand the free shift on $C^*_{\\mathrm{r}} (F_{\\infty})$\nin several examples,\nwe establish notation for them separately.\n\n\\begin{ntn}\\label{N_9214_QFreeFn}\nFor $n \\in {\\mathbb{Z}}_{> 0} \\cup \\{ \\infty \\}$,\nwe let $F_n$ denote the free group on $n$~generators.\nWe let $u_1, u_2, \\ldots, u_n \\in C^*_{\\mathrm{r}} (F_n)$\n(when $n = \\infty$,\nfor $u_1, u_2, \\ldots \\in C^*_{\\mathrm{r}} (F_{\\infty})$)\nbe the ``standard'' unitaries in $C^*_{\\mathrm{r}} (F_n)$,\nobtained as the images of the standard generators of~$F_n$.\nFor $\\zeta = (\\zeta_1, \\zeta_2, \\ldots, \\zeta_n) \\in (S^1)^n$\n(when $n = \\infty$,\nfor $\\zeta = (\\zeta_n)_{n \\in {\\mathbb{Z}}_{> 0}} \\in (S^1)^{{\\mathbb{Z}}_{> 0}}$),\nwe let $\\varphi_{\\zeta} \\in {\\mathrm{Aut}} ( C^*_{\\mathrm{r}} (F_n) )$\nbe the (quasifree) automorphism\ndetermined by $\\varphi_{\\zeta} (u_k) = \\zeta_k u_k$\nfor $k = 1, 2, \\ldots, n$\n(when $n = \\infty$,\nfor $k = 1, 2, \\ldots)$.\nIt is well known that $\\zeta \\mapsto \\varphi_{\\zeta}$ is continuous.\n\\end{ntn}\n\n\\begin{ntn}\\label{N_9214_FreeShift}\nTake the standard generators of the free group $F_{\\infty}$\nto be indexed by~${\\mathbb{Z}}$,\nand for $n \\in {\\mathbb{Z}}$ let $u_n \\in C^*_{\\mathrm{r}} (F_{\\infty})$\nbe the unitary obtained as the image of the corresponding\ngenerator of~$F_{\\infty}$.\nWe denote by $\\sigma$ the free shift on $C^*_{\\mathrm{r}} (F_{\\infty})$,\nthat is, the automorphism\n$\\sigma \\in {\\mathrm{Aut}} ( C^*_{\\mathrm{r}} (F_{\\infty}) )$\ndetermined by $\\sigma (u_n) = u_{n + 1}$\nfor $n \\in {\\mathbb{Z}}$.\n\\end{ntn}\n\n\\begin{exa}\\label{E_9214_ShiftAndUHF}\nSet $X_0 = (S^1)^{{\\mathbb{Z}}}$.\nLet $h_0 \\colon X_0 \\to X_0$\nbe the (forwards) shift,\ndefined for $\\zeta = (\\zeta_n)_{n \\in {\\mathbb{Z}}} \\in X_0$ by\n$h_0 (\\zeta) = (\\zeta_{n - 1})_{n \\in {\\mathbb{Z}}}$.\nLet $D = \\bigotimes_{n \\in {\\mathbb{Z}}_{> 0}} M_2$\nbe the $2^{\\infty}$~UHF algebra.\nChoose a bijection $\\sigma \\colon {\\mathbb{Z}}_{> 0} \\to {\\mathbb{Z}}$,\nand for $\\zeta = (\\zeta_n)_{n \\in {\\mathbb{Z}}} \\in X_0$\ndefine\n\\[\n\\alpha_{\\zeta}^{(0)}\n = \\bigotimes_{n \\in {\\mathbb{Z}}_{> 0}}\n {\\mathrm{Ad}} \\left( \\left( \\begin{matrix}\n 1 & 0 \\\\\n 0 & \\zeta_{\\sigma (n)}\n \\end{matrix} \\right) \\right)\n \\in {\\mathrm{Aut}} (D).\n\\]\nThen $\\zeta \\mapsto \\alpha_{\\zeta}^{(0)}$ is continuous.\n\nWe have $\\mathrm{mdim} (h_0) = 1$ (see Notation~\\ref{N_0108_DimX})\nby Proposition~3.3 of~\\cite{LndWs},\nand we can use $X_0$, $h_0$, $D$, and $\\zeta \\mapsto \\alpha^{(0)}_{\\zeta}$\nin Lemma~\\ref{L_4Y12_GetAuto}.\nHowever, $h_0$ is not minimal.\nWe therefore proceed as follows.\nIdentify $[0, 1]$ with a closed arc in~$S^1$,\nsay via $\\lambda \\mapsto \\exp (\\pi i \\lambda)$.\nUse this identification to identify\n$[0, 1]^{{\\mathbb{Z}}}$ with a closed subset of $X_0 = (S^1)^{{\\mathbb{Z}}}$.\nThis identification is equivariant when both spaces\nare equipped with the shift homeomorphism{s}.\nLet $X \\S X_0$ and $h = h_0 |_{X} \\colon X \\to X$\nbe the minimal subshift\nin~\\cite{GlKr},\nwhich can be taken to have mean dimension\narbitrarily close to~$1$.\nFor $\\zeta \\in X$ let $\\alpha_{\\zeta} = \\alpha_{\\zeta}^{(0)}$.\nLet $\\alpha \\colon {\\mathbb{Z}} \\to {\\mathrm{Aut}} (C (X, D))$\nbe the corresponding action as in Lemma~\\ref{L_4Y12_GetAuto}.\nThen $\\alpha$ lies over the free minimal action of ${\\mathbb{Z}}$ on~$X$\ngenerated by~$h$.\n\nThe crossed product\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is simple\nby Proposition~\\ref{P_4Y15_CPSimple},\nand is nuclear,\nso it is ${\\mathcal{Z}}$-stable by Theorem~\\ref{ATraciallyZStable}.\n\\end{exa}\n\nThe action in Example~\\ref{E_9214_ShiftAndUHF}\ncan also be described as follows.\nRealize $D$ as the algebra of the canonical anticommutation\nrelations on generators $a_k$ for $k \\in {\\mathbb{Z}}_{> 0}$,\nfollowing Section~5.1 of~\\cite{Brtl}.\nThen $\\alpha_{\\zeta}^{(0)}$ is the gauge automorphism\n$\\alpha_{\\zeta}^{(0)} (a_k) = \\zeta_{\\sigma (k)} a_k$.\nIn~\\cite{Brtl}, see Section~5.1 and the proof of Lemma~5.2.\n\nThe next example is a slightly different version,\nwith larger mean dimension.\n\n\\begin{exa}\\label{E_9Z26_ShiftAndUHF_v2}\nLet $X_0$, $h_0$, $D$, and $\\zeta \\mapsto \\alpha_{\\zeta}^{(0)}$\nbe as in Example~\\ref{E_9214_ShiftAndUHF}.\nWe will, however, use the homeomorphism{} $h_0^2$ of~$X_0$,\nwhich is the shift on $(S^1 \\times S^1)^{{\\mathbb{Z}}}$.\nLet $(X, h)$ be the minimal subspace of the shift on $([0, 1]^2)^{{\\mathbb{Z}}}$\nconstructed in Proposition 3.5 of \\cite{LndWs},\nwhich satisfies $\\mathrm{mdim} (h) > 1$ (Notation~\\ref{N_0108_DimX}).\nUse an embedding of $[0, 1]^2$ in~$(S^1)^2$\nto choose an equivariant homeomorphism{} $g$\nfrom $(X, h)$ to an invariant closed subset of $(X_0, h_0^2)$.\nFor $x \\in X$ let $\\alpha_{x} = \\alpha_{g (x)}^{(0)}$.\nLet $\\alpha \\colon {\\mathbb{Z}} \\to {\\mathrm{Aut}} (C (X, D))$\nbe the corresponding action as in Lemma~\\ref{L_4Y12_GetAuto}.\nThe crossed product\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is again\nsimple by Proposition~\\ref{P_4Y15_CPSimple},\nso ${\\mathcal{Z}}$-stable by Theorem~\\ref{ATraciallyZStable}.\n\\end{exa}\n\n\\begin{exa}\\label{E_9214_ShiftAndCStFI_2}\nLet $X$, $h$, $D$, and $\\zeta \\mapsto \\alpha_{\\zeta}$ be as in\nExample~\\ref{E_9214_ShiftAndUHF}.\nDefine $E = D \\otimes C^*_{\\mathrm{r}} (F_{\\infty})$,\nand, following Notation~\\ref{N_9214_QFreeFn},\nfor $\\zeta \\in X$ define\n$\\beta_{\\zeta} = \\alpha_{\\zeta} \\otimes \\varphi_{\\zeta}$.\nApply Lemma~\\ref{L_4Y12_GetAuto}\nto get an action $\\beta \\colon {\\mathbb{Z}} \\to {\\mathrm{Aut}} (C (X, E))$\nwhich lies over the free minimal action of ${\\mathbb{Z}}$ on~$X$\ngenerated by~$h$.\n\nWe claim that\n$C^* \\big( {\\mathbb{Z}}, \\, C (X, E), \\, \\beta \\big)$ is tracially ${\\mathcal{Z}}$-stable\nand has stable rank one.\nTracial ${\\mathcal{Z}}$-stability follows from Theorem~\\ref{ATraciallyZStable}.\nFor stable rank one,\n$E$ is exact,\nso has strict comparison of positive elements\nby Corollary 4.6 of~\\cite{Rdm7}.\nChoose any one point subset $Y \\S X$.\nThen $C^* \\big( {\\mathbb{Z}}, \\, C (X, E), \\, \\beta \\big)_Y$\n(see Definition~\\ref{D_4Y12_OrbSubalg})\nis ${\\mathcal{Z}}$-stable\nby Corollary~\\ref{EStableGivesZStable}.\nThis algebra\nis simple by Proposition~\\ref{P_4Y15_CPSimple},\nso by Theorem 6.7 of~\\cite{Rdm7} it has stable rank one.\nThe algebra $C^* \\big( {\\mathbb{Z}}, \\, C (X, E), \\, \\beta \\big)_Y$\nis centrally large\nin $C^* \\big( {\\mathbb{Z}}, \\, C (X, E), \\, \\beta \\big)$\nby\nCorollary \\ref{AyCentLargeConds}(\\ref{9212_AyCentLargeConds_StrCmp}),\nso Theorem 6.3 of~\\cite{ArPh} implies\nthat $C^* \\big( {\\mathbb{Z}}, \\, C (X, E), \\, \\beta \\big)$ has stable rank one.\n\\end{exa}\n\nIn Example~\\ref{E_9214_ShiftAndCStFI_2}, we don't know whether\n$C^* \\big( {\\mathbb{Z}}, \\, C (X, E), \\, \\beta \\big)$ is ${\\mathcal{Z}}$-stable.\nWe also don't know whether tracial ${\\mathcal{Z}}$-stability implies\nstable rank one, although this is expected to be true.\n\nThere is nothing special about the specific formulas\nfor $x \\mapsto \\alpha_x$\nin Example~\\ref{E_9214_ShiftAndUHF}\nand Example~\\ref{E_9Z26_ShiftAndUHF_v2},\nand $x \\mapsto \\beta_x$ in Example~\\ref{E_9214_ShiftAndCStFI_2}.\nThey were chosen merely to show that interesting examples exist.\n\nExample~\\ref{E_9214_ShiftAndUHF},\nExample~\\ref{E_9Z26_ShiftAndUHF_v2},\nand Example~\\ref{E_9214_ShiftAndCStFI_2}\nwere constructed so that the homeomorphism~$h$\ndoes not have mean dimension zero.\nIf $X$ is finite dimensional,\nthen one can get all we do by using known results\nfor crossed products by actions with finite Rokhlin dimension\nwith commuting towers.\nWith finite dimensional~$X$, one even gets ${\\mathcal{Z}}$-stability in\nexamples like Example~\\ref{E_9214_ShiftAndUHF},\nExample~\\ref{E_9Z26_ShiftAndUHF_v2},\nand Example~\\ref{E_9214_ShiftAndCStFI_2},\nby Theorem~5.8 of~\\cite{HWZ}.\nHowever,\nwe know of no results which apply to examples of this type\nwhen $X$ is infinite dimensional{} and $h$ has mean dimension zero.\n\nThe next most obvious choice for a minimal homeomorphism{} of an infinite dimensional{} space~$X$\nseems to be as follows.\nTake $X = (S^1)^{{\\mathbb{Z}}_{> 0}}$,\nfix $\\theta = (\\theta_n)_{n \\in {\\mathbb{Z}}_{> 0}} \\in {\\mathbb{R}}^{{\\mathbb{Z}}_{> 0}}$,\nand define $h \\colon X \\to X$ by\n$h \\bigl( (\\zeta_{n})_{n \\in {\\mathbb{Z}}_{> 0}} \\bigr)\n = \\bigl( e^{2 \\pi i \\theta_n} \\zeta_{n} \\bigr)_{n \\in {\\mathbb{Z}}_{> 0}}$.\nIf $1, \\theta_1, \\theta_2, \\ldots$ are linearly independent over~${\\mathbb{Q}}$,\nthen $h$ is minimal\nby Proposition~\\ref{P_4X04_ProdRotation}.\nHowever, by considering the action in just one coordinate,\none sees that, regardless of $\\zeta \\mapsto \\alpha_{\\zeta}$,\nthe action has finite Rokhlin dimension with commuting towers.\nBy Theorem 6.2 of~\\cite{HWZ}, irrational rotations\nhave finite Rokhlin dimension with commuting towers.\nRemark 6.3 of~\\cite{HWZ}, according to which the result\nextends to any homeomorphism which has an irrational\nrotation as a factor,\napplies equally well to any automorphism of $C (S^1) \\otimes B$,\nfor any unital C*-algebra~$B$,\nwhich lies over an irrational rotation on~$S^1$\nin the sense of Definition~\\ref{D_4Y12_OverX}.\nIf we start with $C \\bigl( (S^1)^{{\\mathbb{Z}}_{> 0}}, \\, D \\bigr)$,\ntake $B$ above to be $C \\bigl( (S^1)^{{\\mathbb{Z}}_{> 0} \\setminus \\{ 1 \\}}, \\, D \\bigr)$.\nTherefore Theorem~5.8 of~\\cite{HWZ} applies,\nand our results give nothing new.\n\nWe now give some examples of the second type discussed\nin the introduction to this section.\nFor easy reference,\nwe recall a result on the stable rank of reduced free products.\nMany reduced free products have stable rank~$1$.\nThe following result is from~\\cite{DHR}.\n(It is not affected by the correction~\\cite{DHRc}.)\nReduced free products of unital C*-algebra{s}\nin~\\cite{DHR}\nare implicitly taken to be amalgamated over~${\\mathbb{C}}$;\nsee Section~2.2 of~\\cite{DHR}.\n\n\\begin{prp}[Corollary~3.9 of~\\cite{DHR}]\\label{P_5416_TsrFrGp}\nLet $G$ and $H$ be discrete groups\nwith ${\\mathrm{card}} (G) \\geq 2$\nand ${\\mathrm{card}} (H) \\geq 3$.\nThen $C^*_{\\mathrm{r}} (G \\star H)$ has stable rank~$1$.\n\\end{prp}\n\n\\begin{exa}\\label{E_5416_CSrFGp}\nLet $X$ be the Cantor set\nand let $h$ be an arbitrary minimal\nhomeomorphism of~$X$.\nLet $\\sigma \\in {\\mathrm{Aut}} ( C^*_{\\mathrm{r}} (F_{\\infty}) )$\nbe as in Notation~\\ref{N_9214_FreeShift}.\nLet\n$\\alpha \\in {\\mathrm{Aut}} \\big( C (X) \\otimes C^*_{\\mathrm{r}} (F_{\\infty}) \\big)$\nbe the tensor product of the automorphism\n$f \\mapsto f \\circ h^{-1}$ of $C (X)$ and~$\\sigma$.\nThen\n$C^* \\big( {\\mathbb{Z}}, \\, C (X) \\otimes C^*_{\\mathrm{r}} (F_{\\infty}),\n \\, \\alpha \\big)$\nis simple by Proposition~\\ref{P_4Y15_CPSimple}.\n\nWe claim that\n$C^* \\big( {\\mathbb{Z}}, \\, C (X) \\otimes C^*_{\\mathrm{r}} (F_{\\infty}),\n \\, \\alpha \\big)$\nhas stable rank~$1$.\n\nTo prove the claim,\nwe apply Theorem~\\ref{P_9Z26_XisCantor} with\n$D = C^*_{\\mathrm{r}} (F_{\\infty})$ and\n$\\alpha$ as given.\nUse Theorems 9.2.6 and 9.2.7 of \\cite{PhNotes}\nto see that $D$ is simple and has a tracial state.\nBy Proposition~\\ref{P_5416_TsrFrGp}, we have ${\\mathrm{tsr}} (D) = 1$.\nBy Proposition 6.3.2 of \\cite{Rbt},\n$D$ has strict comparison of positive elements.\nBy construction, $\\alpha$ lies over~$h$.\nNow apply\nTheorem~\\ref{P_9Z26_XisCantor}(\\ref{9217_GetTsr1_StrCmp}).\n\\end{exa}\n\nApparently no previously known results give\nanything about the crossed product in this example.\nIn particular, as discussed after Problem~\\ref{Pb_9922_RokhlinTsr},\nknowing that the action has the Rokhlin property\ndoesn't seem to help.\n\nThere are many other automorphisms\nof $C^*_{\\mathrm{r}} (F_{\\infty})$\nwhich could be used in place of~$\\sigma$.\nFor example,\none could take a quasifree automorphism,\nas in Notation~\\ref{N_9214_QFreeFn}.\nHere is a more interesting version.\n\n\\begin{exa}\\label{E_5416_DenjFst}\nFix any $n \\in \\set{ 2, 3, \\ldots, \\infty }$\nand take $D = C^*_{\\mathrm{r}} (F_{n})$.\nAdopt Notation~\\ref{N_9214_QFreeFn}.\nLet $h \\colon X \\to X$ be a restricted Denjoy homeomorphism{}\nwith rotation number $\\theta \\in {\\mathbb{R}} \\setminus {\\mathbb{Q}}$,\nas in Definition~\\ref{D_9X04_RestrDj}.\nLet $\\zeta \\colon X \\to S^1$ be the continuous{} surjective map\nwith $\\zeta (h (x)) = e^{2 \\pi i \\theta} \\zeta (x)$\nof Lemma~\\ref{L_9X04_DjSemiCj}\n(gotten from Corollary~3.2 and Proposition~3.4 of~\\cite{PtSmdSk}).\nFollowing Notation~\\ref{N_9214_QFreeFn} for the generators\nof~$D$,\nfor $x \\in X$\nlet $\\alpha_x \\in {\\mathrm{Aut}} ( D )$\nbe determined\nby $\\alpha_x (u_k) = \\zeta (x) u_k$\nfor $k = 1, 2, \\ldots, n$ (or $k \\in {\\mathbb{Z}}_{> 0}$ if $n = \\infty$).\nApply Lemma~\\ref{L_4Y12_GetAuto}\nto get an action\n$\\alpha \\colon {\\mathbb{Z}} \\to {\\mathrm{Aut}} \\big( C (X, D ) \\big)$,\nwhich we can think of as\na kind of noncommutative Furstenberg transformation.\nThe algebra $C^* \\big( {\\mathbb{Z}}, \\, C (X, D ), \\, \\alpha \\big)$\nis simple by Proposition~\\ref{P_4Y15_CPSimple}.\n\nWe claim that\n$C^* \\big( {\\mathbb{Z}}, \\, C (X, D ), \\, \\alpha \\big)$\nhas stable rank one.\nTo prove the claim,\nuse Theorems 9.2.6 and 9.2.7 of \\cite{PhNotes}\nto see that $D$ is simple and has a tracial state.\nBy Proposition~\\ref{P_5416_TsrFrGp}, we have ${\\mathrm{tsr}} (D) = 1$.\nNow apply\nTheorem~\\ref{P_9Z26_XisCantor}(\\ref{9217_GetTsr1_CptGen}).\n\\end{exa}\n\nAgain, apparently no previously known results give\nanything about the crossed product.\nSince $C^*_{\\mathrm{r}} (F_{n})$ is not ${\\mathcal{Z}}$-stable,\nknowing that the action has the Rokhlin property\ndoesn't seem to help.\n\nWe can generalize Example~\\ref{E_5416_DenjFst} as follows.\n\n\\begin{exa}\\label{E_9217_MultltDenj}\nFix $n \\in \\{ 2, 3, \\ldots, \\infty \\}$.\nFix $\\theta_1, \\theta_2, \\ldots, \\theta_n \\in {\\mathbb{R}}$\n(or $\\theta_1, \\theta_2, \\ldots \\in {\\mathbb{R}}$)\nsuch that $1, \\theta_1, \\theta_2, \\ldots, \\theta_n$\n(or $1, \\theta_1, \\theta_2, \\ldots$)\nare linearly independent over~${\\mathbb{Q}}$.\nFor $k = 1, 2, \\ldots, n$,\nlet $h \\colon X_k \\to X_k$ be a restricted Denjoy homeomorphism{}\nwith rotation number $\\theta_k \\in {\\mathbb{R}} \\setminus {\\mathbb{Q}}$,\nas in Definition~\\ref{D_9X04_RestrDj}.\nTake $X = \\prod_{k = 1}^n X_k$,\nand let $h \\colon X \\to X$\nact as $h_k$ on the $k$-th factor.\nThen~$h$ is minimal by Proposition~\\ref{P_4X04_ProdRstDj}.\nDefine $\\zeta_k \\colon X_k \\to S^1$\nanalogously to the definition of $\\zeta$\nin Example~\\ref{E_5416_DenjFst},\nand take $\\alpha_x (u_k) = \\zeta_k (x_k) u_k$.\nApply Lemma~\\ref{L_4Y12_GetAuto}\nto get an action\n$\\alpha \\colon\n {\\mathbb{Z}} \\to {\\mathrm{Aut}} \\big( C (X, \\, C^*_{\\mathrm{r}} (F_{n}) ) \\big)$.\nThen\n$C^* \\big( {\\mathbb{Z}}, \\, C (X, \\, C^*_{\\mathrm{r}} (F_{n}) ),\n \\, \\alpha \\big)$\nis simple and has stable rank one\nfor the same reasons as in Example~\\ref{E_5416_DenjFst}.\n\\end{exa}\n\n\\begin{exa}\\label{R_5416_More}\nLet $X$ be the Cantor set\nand let $h$ be an arbitrary minimal\nhomeomorphism of~$X$.\nChoose any decomposition $X = X_1 \\amalg X_2$\nof $X$ as the disjoint union of two nonempty closed subsets.\nSet $D = C^*_{\\mathrm{r}} (F_{\\infty})$,\nwith the generators of $F_{\\infty}$ indexed by~${\\mathbb{Z}}$.\nFix any $\\zeta = (\\zeta_n)_{n \\in {\\mathbb{Z}}} \\in (S^1)^{{\\mathbb{Z}}_{> 0}}$.\nFor $x \\in X_1$ take\n$\\alpha_{x}$ to be the automorphism\n$\\varphi_{\\zeta}$ of Notation~\\ref{N_9214_QFreeFn},\nexcept using ${\\mathbb{Z}}$ in place of ${\\mathbb{Z}}_{> 0}$,\nand for $x \\in X_2$ take\n$\\alpha_{x}$ to be the automorphism\n$\\sigma$ of Notation~\\ref{N_9214_FreeShift}.\nApply Lemma~\\ref{L_4Y12_GetAuto}\nto get an action\n$\\alpha \\colon {\\mathbb{Z}} \\to {\\mathrm{Aut}} ( C (X, D) )$.\nThe algebra $C^* ({\\mathbb{Z}}, \\, C (X) \\otimes D, \\, \\alpha)$\nis simple by Proposition~\\ref{P_4Y15_CPSimple}.\nThe algebra $D$ has strict comparison of positive elements by\nProposition 6.3.2 of~\\cite{Rbt},\nso $C^* ({\\mathbb{Z}}, \\, C (X) \\otimes D, \\, \\alpha)$ has\nstable rank one\nby Theorem~\\ref{P_9Z26_XisCantor}(\\ref{9217_GetTsr1_StrCmp}).\n\\end{exa}\n\nExample~\\ref{R_5416_More} admits many variations.\nHere are several.\n\n\\begin{exa}\\label{E_9217_AltDenj}\nLet $h \\colon X \\to X$ be a restricted Denjoy homeomorphism,\nas in Example~\\ref{E_5416_DenjFst},\nand let $\\zeta \\colon X \\to S^1$ be as there.\nChoose any decomposition $X = X_1 \\amalg X_2$\nof $X$ as the disjoint union of two nonempty closed subsets.\nFollowing Notation~\\ref{N_9214_QFreeFn} for the generators\nof $C^*_{\\mathrm{r}} (F_{2})$,\nfor $x \\in X_1$\nlet $\\alpha_x \\in {\\mathrm{Aut}} ( C^*_{\\mathrm{r}} (F_{n}) )$\nbe determined\nby $\\alpha_x (u_1) = \\zeta (x) u_1$ and $\\alpha_x (u_2) = u_2$,\nand for $x \\in X_2$\nlet $\\alpha_x \\in {\\mathrm{Aut}} ( C^*_{\\mathrm{r}} (F_{n}) )$\nbe determined\nby $\\alpha_x (u_1) = u_1$ and $\\alpha_x (u_2) = \\zeta (x) u_2$.\nThe algebra\n$C^* \\big( {\\mathbb{Z}}, \\, C (X)\n \\otimes C^*_{\\mathrm{r}} (F_{2}), \\, \\alpha \\big)$\nis simple by Proposition~\\ref{P_4Y15_CPSimple}.\n\nWe claim that\nthis algebra\nhas stable rank one.\nTo prove the claim,\nuse Theorems 9.2.6 and 9.2.7 of \\cite{PhNotes}\nto see that $D$ is simple and has a tracial state.\nBy Proposition~\\ref{P_5416_TsrFrGp}, we have ${\\mathrm{tsr}} (D) = 1$.\nNow apply\nTheorem~\\ref{P_9Z26_XisCantor}(\\ref{9217_GetTsr1_CptGen}).\n\\end{exa}\n\n\\begin{exa}\\label{E_9217_DenjAndSw}\nIn Example~\\ref{E_9217_AltDenj},\nreplace the definition of $\\alpha_x$\nfor $x \\in X_2$ with\n$\\alpha_x (u_1) = u_2$ and $\\alpha_x (u_2) = u_1$.\nProposition~\\ref{P_4Y15_CPSimple}\nand Theorem~\\ref{P_9Z26_XisCantor}(\\ref{9217_GetTsr1_CptGen})\nstill apply,\nso\n$C^* \\big( {\\mathbb{Z}}, \\, C (X)\n \\otimes C^*_{\\mathrm{r}} (F_{2}), \\, \\alpha \\big)$\nagain is simple and has stable rank one.\n\\end{exa}\n\nWe now give examples in which we also get real rank zero.\nThe following lemma will be used for them.\n\n\\begin{lem}\\label{L_0916_MakeSimple}\nLet $M$ be factor of type $\\mathrm{II}_1$,\nlet $\\alpha \\colon G \\to {\\mathrm{Aut}} (M)$\nbe an action of a countable group~$G$ on~$M$,\nand let $P \\S M$ be a separable C*-subalgebra.\nThen there exists a simple separable unital C*-subalgebra\n$D \\subset M$ which contains~$P$,\nis invariant under ${\\overline{\\sigma}}$,\nhas a unique tracial state\n(the restriction to $D$ of the unique tracial state on $M$),\nhas real rank zero and stable rank one,\nand such that the order on projections over~$D$\nis determined by traces.\n\\end{lem}\n\n\\begin{proof}\nThe proof is the same as that of Proposition~3.1 of~\\cite{Ph57},\nexcept for $G$-invariance.\nWe construct by induction on $n \\in {\\mathbb{Z}}_{\\geq 0}$ separable unital subalgebras\n$A_n, B_n, C_n, D_n, E_n, F_n \\S M$ with\n\\[\nP \\S A_0 \\S B_0 \\S C_0 \\S D_0 \\S E_0 \\S F_0 \\S \\cdots\n \\S A_n \\S B_n \\S C_n \\S D_n \\S E_n \\S F_n \\S \\cdots\n\\]\nsuch that $A_n$, $B_n$, $C_n$, $D_n$, and $E_n$ have the\nproperties in the proof of Proposition~3.1 of~\\cite{Ph57}\n($A_n$ is simple, etc.),\nand such that $F_n$ is $G$-invariant.\nIn the induction step, after constructing $E_n$,\nwe take $F_n$ to be the C*-subalgebra generated by\n$\\bigcup_{g \\in G} \\alpha_g (E_n)$.\nThen ${\\overline{\\bigcup_{n = 0}^{\\infty} F_n}}$ is $G$-invariant.\nThe proof in~\\cite{Ph57} now gives the desired conclusion,\nexcept that the conclusion is that $K_0 (D) \\to K_0 (M)$\nis an order isomorphism onto its range\ninstead of that the order on projections over~$D$\nis determined by traces.\nBut the conclusion we get implies the conclusion we want\nif $A$ has cancellation,\nand stable rank one implies cancellation\nby Proposition 6.4.1 and Proposition 6.5.1 of~\\cite{Blkd3}.\n\\end{proof}\n\n\\begin{exa}\\label{E_5416_DenjRR0}\nLet $X$ be the Cantor set\nand let $h$ be an arbitrary minimal\nhomeomorphism of~$X$.\nLet $\\sigma \\in {\\mathrm{Aut}} ( C^*_{\\mathrm{r}} (F_{\\infty}) )$\nbe as in Notation~\\ref{N_9214_FreeShift}.\n\nWe regard $C^*_{\\mathrm{r}} (F_{\\infty})$\nas a subalgebra of the group\nvon Neumann algebra $W^* (F_{\\infty})$\nin the usual way,\nand we let ${\\overline{\\sigma}} \\in {\\mathrm{Aut}} (W^* (F_{\\infty}))$\nbe the von Neumann algebra automorphism\nwhich shifts the generators of $F_{\\infty}$\nin the same way that $\\sigma$ does.\nThus ${\\overline{\\sigma}} |_{C^*_{\\mathrm{r}} (F_{\\infty})} = \\sigma$.\nChoose a ${\\overline{\\sigma}}$-invariant\nsubalgebra $D \\subset W^* (F_{\\infty})$\nwhich contains $C^*_{\\mathrm{r}} (F_{\\infty})$\nas in Lemma~\\ref{L_0916_MakeSimple},\nwith the properties there.\nSet $\\gamma = {\\overline{\\sigma}} |_D \\in {\\mathrm{Aut}} (D)$.\nThen define $\\alpha \\in {\\mathrm{Aut}} ( C (X) \\otimes D )$\nin the same was as in Example~\\ref{E_5416_CSrFGp},\nusing $\\gamma$ in place of~$\\sigma$.\nThe algebra $C^* \\big( {\\mathbb{Z}}, \\, C (X) \\otimes D, \\, \\alpha \\big)$\nis simple by Proposition~\\ref{P_4Y15_CPSimple}.\nIt has stable rank~$1$ and real rank zero\nby Theorem \\ref{P_9Z26_XisCantor}(\\ref{9217_GetTsr1_SP}).\n\nThe algebra $D$ is not nuclear\nbecause the Gelfand-Naimark-Segal representation\nfrom its tracial state gives a nonhyperfinite factor.\n\\end{exa}\n\nIt is perhaps interesting to point out that every quasitrace\non the algebra~$D$ in Example~\\ref{E_5416_DenjRR0} is a trace.\nThis is a consequence of the following lemma.\n\n\\begin{lem}\\label{L_9929_UniqQTr}\nLet $A$ be a unital C*-algebra{} with real rank zero.\nSuppose that $\\tau$ is a quasitrace on~$A$ and\nthat whenever $p, q \\in M_{\\infty} (A)$ are projection{s}\nwith $\\tau (p) < \\tau (q)$,\nthen $p \\precsim q$.\nThen $\\tau$ is the only quasitrace on~$A$.\n\\end{lem}\n\n\\begin{proof}\nSuppose the conclusion is false,\nand let $\\sigma$ be some other quasitrace on~$A$.\nBy definition, for any quasitrace~$\\rho$ on~$A$\nand any $a, b \\in A_{\\mathrm{sa}}$,\nwe have $\\rho (a + i b) = \\rho (a) + i \\rho (b)$.\nTherefore there is $a \\in A_{\\mathrm{sa}}$\nsuch that $\\sigma (a) \\neq \\tau (a)$.\nSince quasitraces are continuous,\nit follows from real rank zero that there is $a \\in A_{\\mathrm{sa}}$\nsuch that $a$ has finite spectrum and $\\sigma (a) \\neq \\tau (a)$.\nSince quasitraces are linear on commutative C*-subalgebras,\nthere is a projection{} $p \\in A$ such that $\\sigma (p) \\neq \\tau (p)$.\n\nSuppose $\\sigma (p) < \\tau (p)$.\nChoose $m, n \\in {\\mathbb{Z}}_{\\geq 0}$ such that $m \\sigma (p) < n < m \\tau (p)$.\nDefine projection{s} $e, f \\in M_{\\infty} (A)$ by $e = 1_{M_m} \\otimes p$ and\n$f = 1_{M_n} \\otimes 1_A$.\nThen $\\tau (f) < \\tau (e)$, so $f \\precsim e$.\nBut $\\sigma (f) > \\sigma (e)$, a contradiction.\nSimilarly, $\\sigma (p) > \\tau (p)$ is also impossible.\nThis contradiction shows that $\\sigma$ does not exist.\n\\end{proof}\n\n\\begin{exa}\\label{E_9217_BiggerCstFn}\nLet $n \\in \\set{ 2, 3, \\ldots }$.\nFollowing Notation~\\ref{N_9214_QFreeFn} for the generators\nof $C^*_{\\mathrm{r}} (F_{n})$,\nfor $\\rho$ in the symmetric group~$S_n$\nlet $\\psi_{\\rho} \\in {\\mathrm{Aut}} ( C^*_{\\mathrm{r}} (F_{n}) )$\nbe the automorphism determined\nby $\\psi_{\\rho} (u_k) = u_{\\rho^{-1} (k)}$ for $k = 1, 2, \\ldots, n$,\nand let ${\\overline{\\psi_{\\rho} }}$ be the corresponding automorphism\nof the group von Neumann algebra $W^* (F_{n})$.\n\nBy Theorems 9.2.6 and 9.2.7 of \\cite{PhNotes},\nthe algebra $C^*_{\\mathrm{r}} (F_{n})$ is simple\nand has a unique tracial state.\nUse Lemma~\\ref{L_0916_MakeSimple}\nto find a simple separable unital C*-algebra{}\n$D \\subset W^* (F_{n})$ which contains $C^*_{\\mathrm{r}} (F_{n})$,\nis invariant under all the automorphisms\n${\\overline{\\psi_{\\rho} }}$ for $\\rho \\in S_n$,\nand has the other properties given in Lemma~\\ref{L_0916_MakeSimple}.\n\nLet $X$ be the Cantor set, and let $h \\colon X \\to X$ be any minimal homeomorphism.\nChoose any decomposition $X = \\coprod_{\\rho \\in S_n} X_{\\rho}$\nof $X$ as the disjoint union of nonempty closed subsets.\nFor $x \\in X_{\\rho}$ let $\\alpha_x = {\\overline{\\psi_{\\rho} }} |_D$.\nLet $\\alpha \\colon {\\mathbb{Z}} \\to {\\mathrm{Aut}} (C (X, D))$\nbe the corresponding action as in Lemma~\\ref{L_4Y12_GetAuto}.\n\nThe algebra $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nis simple by Proposition~\\ref{P_4Y15_CPSimple}.\nIt has stable rank one and real rank zero by\nTheorem~\\ref{P_9Z26_XisCantor}(\\ref{9217_GetTsr1_CptGen}).\nThe algebra $D$ is not nuclear\nfor the same reason as in Example~\\ref{E_5416_DenjRR0}.\n\\end{exa}\n\nThe next example is of the third type discussed\nin the introduction to this section.\n\n\\begin{exa}\\label{E_9217_GenPType}\nLet $(\\pi_1, \\pi_2, \\ldots)$\nbe a sequence of unital finite dimensional{} representations\nof $C^* (F_2)$\nsuch that, for every $n \\in {\\mathbb{Z}}_{> 0}$,\nthe representation\n$\\bigoplus_{k = n}^{\\infty} \\pi_k$ is faithful.\nFor $n \\in {\\mathbb{Z}}_{\\geq 0}$,\nlet $d (n)$ be the dimension of $\\pi_n$,\nand define\n$l (n) = d (n) + 4$\nand $r (n) = \\prod_{k = 1}^n l (k)$.\nLet $u_1, u_2 \\in C^* (F_2)$\nbe the ``standard'' unitaries in $C^* (F_2)$,\nobtained as the images of the standard generators of~$F_2$,\nas in Notation~\\ref{N_9214_FreeShift}\nexcept that we are now using the full C*-algebra{} instead\nof the reduced C*-algebra.\nLet $\\gamma_1, \\gamma_2, \\gamma_3 \\in {\\mathrm{Aut}} (C^* (F_2))$\nbe the automorphisms determined by\n\\[\n\\gamma_1 (u_1) = u_1^{-1}\n\\qquad {\\mbox{and}} \\qquad\n\\gamma_1 (u_2) = u_2,\n\\]\n\\[\n\\gamma_2 (u_1) = u_1\n\\qquad {\\mbox{and}} \\qquad\n\\gamma_2 (u_2) = u_2^{-1},\n\\]\nand\n\\[\n\\gamma_3 (u_1) = u_1^{-1}\n\\qquad {\\mbox{and}} \\qquad\n\\gamma_3 (u_2) = u_2^{-1}.\n\\]\nFor $n \\in {\\mathbb{Z}}_{> 0}$\ndefine\n$D_n = M_{r (n)} \\otimes C^* (F_2)$,\nwhich we identify as\n\\[\nM_{l (n)} \\otimes M_{l (n - 1)} \\otimes \\cdots\n \\otimes M_{l (1)} \\otimes C^* (F_2),\n\\]\nand define\n$\\gamma_{n, n - 1} \\colon D_{n - 1} \\to D_n$\nby,\nfor $a \\in D_{n - 1} = M_{r (n - 1)} \\otimes C^* (F_2)$,\n\\begin{align*}\n\\gamma_{n, n - 1} (a)\n& = {\\mathrm{diag}} \\bigl( a, \\, ({\\mathrm{id}}_{M_{r (n - 1)}} \\otimes \\varphi_1 ) (a),\n \\, ({\\mathrm{id}}_{M_{r (n - 1)}} \\otimes \\varphi_2 ) (a),\n\\\\\n& \\hspace*{3em} {\\mbox{}}\n \\, ({\\mathrm{id}}_{M_{r (n - 1)}} \\otimes \\varphi_3 ) (a),\n \\, (\\pi_n \\otimes {\\mathrm{id}}_{M_{r (n - 1)}} (a) \\otimes 1\n \\bigr)\n\\\\\n& \\in M_{l (n)} \\otimes D_{n - 1} = D_n.\n\\end{align*}\nLet $D$ be the direct limit of the resulting direct system.\nIt follows by methods of~\\cite{Ddlt}\nthat $D$ is simple, separable, and has tracial rank zero.\nIn particular, $D$ has stable rank one\nand real rank zero by Theorem 3.4 of~\\cite{HLinTrAF},\nand the order on projections over $D$ is determined by traces\nby Theorem 6.8 of~\\cite{HLinTTR}.\nMoreover, using the known result for $K_* (C^* (F_n))$\n(see~\\cite{CuFree}),\none gets $K_1 (D) = 0$.\n\nFor $\\zeta \\in S^1$ define inductively\nunitaries $u_n (\\zeta) \\in D_n$ as follows.\nSet $u_0 (\\zeta) = 1$,\nand, given $u_{n - 1} (\\zeta) \\in D_{n - 1}$,\nset\n\\[\nu_n (\\zeta)\n = \\bigl( {\\mathrm{diag}} \\bigl( \\zeta, 1, 1, \\ldots, 1 \\bigr)\n \\otimes 1_{D_{n - 1}} \\bigr)\n \\gamma_{n, n - 1} ( u_{n - 1} (\\zeta) )\n \\in M_{l (n)} \\otimes D_{n - 1} = D_n.\n\\]\nThen the actions\n$\\zeta \\mapsto {\\mathrm{Ad}} (u_n (\\zeta)) \\in {\\mathrm{Aut}} (D_n)$\nof $S^1$ on~$D_n$\nare compatible with the direct system,\nand so yield a continuous{} map\n(in fact, an action)\n$\\zeta \\mapsto \\alpha_{\\zeta} \\colon S^1 \\to {\\mathrm{Aut}} (D)$.\nLet $h \\colon S^1 \\to S^1$ be an irrational rotation.\nApply Lemma~\\ref{L_4Y12_GetAuto}\nto get an action\n$\\alpha \\colon {\\mathbb{Z}} \\to {\\mathrm{Aut}} ( C (S^1, D) )$.\n\nThe algebra $C^* \\big( {\\mathbb{Z}}, \\, C (S^1, D), \\, \\alpha \\big)$\nis simple by Proposition~\\ref{P_4Y15_CPSimple}\nand has stable rank one\nby\nTheorem \\ref{T_9222_Tsr1AndRR0_CrPrd}(\\ref{9927_AyCentLargeConds_SP}).\n\\end{exa}\n\nMethods of~\\cite{NiuWng} will probably show that\nthe algebra $D$ in Example~\\ref{E_9217_GenPType} is not ${\\mathcal{Z}}$-stable\n(see Question~\\ref{Q_9217_GenPType_IsZStab}),\nalthough it definitely is tracially ${\\mathcal{Z}}$-stable.\n\nFinally, we give purely infinite examples.\n\n\\begin{exa}\\label{E_9214_ShiftAndOI}\nLet $(X, h)$ be the minimal subshift\n(of the shift on $(S^1)^{{\\mathbb{Z}}}$) with nonzero mean dimension\nused in Example~\\ref{E_9214_ShiftAndUHF}.\nLet $D = {\\mathcal{O}}_{\\infty}$,\nbut with standard generators indexed by~${\\mathbb{Z}}$,\nsay $s_k$ for $k \\in {\\mathbb{Z}}$.\nLet $\\alpha_{\\zeta}$ be the gauge automorphism,\ngiven by $\\alpha_{\\zeta} (s_k) = \\zeta_k s_k$.\nLet $\\alpha \\colon {\\mathbb{Z}} \\to {\\mathrm{Aut}} (C (X, D))$\nbe the corresponding action as in Lemma~\\ref{L_4Y12_GetAuto}.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is\npurely infinite and simple by Theorem~\\ref{T_0101_PICase}.\nSince $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nis nuclear,\nit is necessarily ${\\mathcal{O}}_{\\infty}$-stable\nby Theorem 3.15 of \\cite{KiPh}.\n\\end{exa}\n\nIn Example~\\ref{E_9214_ShiftAndOI},\nTheorems 4.1 and~5.8 of~\\cite{HWZ} don't apply,\nsince there is no reason to think that $\\alpha$\nhas finite Rokhlin dimension with commuting towers.\n\n\\begin{exa}\\label{E_9214_OnPINotZStab}\nLet $D$ be the reduced free product\n$D = (M_2 \\otimes M_2) \\star_{\\mathrm{r}} C ([0, 1])$,\ntaken with respect to the Lebesgue measure state on $C ([0, 1])$\nand the state on $M_2 \\otimes M_2$\ngiven by tensor product of the usual tracial state ${\\mathrm{tr}}$\nwith the state\n$\\rho (x) = {\\operatorname{tr}} \\big(\n {\\mathrm{diag}} \\big( \\frac{1}{3}, \\frac{2}{3} \\big) x \\big)$\non $M_2$.\nIt is shown in Example 5.8 of~\\cite{AmGlJmPh}\nthat $D$ is purely infinite and simple\nbut not ${\\mathcal{Z}}$-stable.\n\nTake $X = (S^1)^4$,\nand for $\\zeta = (\\zeta_1, \\zeta_2, \\zeta_3, \\zeta_4) \\in X$\ntake $\\alpha_{\\zeta}$\nto be the free product automorphism\nwhich is given by\n\\[\n{\\mathrm{Ad}} \\left( \\left( \\begin{matrix}\n \\zeta_1 & 0 \\\\\n 0 & \\zeta_2\n \\end{matrix} \\right) \\right)\n\\otimes {\\mathrm{Ad}} \\left( \\left( \\begin{matrix}\n \\zeta_3 & 0 \\\\\n 0 & \\zeta_4\n \\end{matrix} \\right) \\right)\n\\]\non $M_2 \\otimes M_2$\nand is trivial on $C ([0, 1])$.\nChoose $\\theta_1, \\theta_2, \\theta_3, \\theta_4 \\in {\\mathbb{R}}$\nsuch that $1, \\theta_1, \\theta_2, \\theta_3, \\theta_4$\nare linearly independent over~${\\mathbb{Q}}$.\nTake $h \\colon X \\to X$\nto be\n\\[\n(\\zeta_1, \\zeta_2, \\zeta_3, \\zeta_4)\n \\mapsto \\bigl( e^{2 \\pi i \\theta_1} \\zeta_1, \\,\n e^{2 \\pi i \\theta_2} \\zeta_2, \\,\n e^{2 \\pi i \\theta_3} \\zeta_3, \\,\n e^{2 \\pi i \\theta_4} \\zeta_4 \\bigr).\n\\]\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is\npurely infinite and simple by Theorem~\\ref{T_0101_PICase}.\n\\end{exa}\n\nThe action in Example~\\ref{E_9214_OnPINotZStab}\nhas finite Rokhlin dimension with commuting towers.\nHowever, we don't know any theorem on pure infiniteness\nfor crossed products by such actions\nwhen the original algebra is not ${\\mathcal{O}}_{\\infty}$-stable.\nWe address this in Question~\\ref{Pb_0111_FinRDim_PI} below.\nOur result also does not imply that the crossed product\nis ${\\mathcal{O}}_{\\infty}$-stable, or even ${\\mathcal{Z}}$-stable,\nalthough it seems plausible that it might be.\n\nIn the next two examples, $D$ is again not ${\\mathcal{Z}}$-stable.\nAlso, $X$ isn't finite dimensional, and $h$ doesn't even have mean dimension zero.\nSo a positive answer to Question~\\ref{Pb_0111_FinRDim_PI}\npresumably would not help.\n\n\\begin{exa}\\label{E_9214_OnPINotZStab_2}\nLet $D$ be as in Example~\\ref{E_9214_OnPINotZStab},\nand let $(X, h)$ be as in Example~\\ref{E_9214_ShiftAndUHF}\n(and reused in Example~\\ref{E_9214_ShiftAndOI}).\nFor $\\zeta = (\\zeta_n)_{n \\in {\\mathbb{Z}}} \\in X$\ntake $\\alpha_{\\zeta}$ to be the free product automorphism\nwhich is given by\n\\[\n{\\mathrm{Ad}} \\left( \\left( \\begin{matrix}\n \\zeta_1 & 0 \\\\\n 0 & \\zeta_2\n \\end{matrix} \\right) \\right)\n\\otimes {\\mathrm{Ad}} \\left( \\left( \\begin{matrix}\n \\zeta_3 & 0 \\\\\n 0 & \\zeta_4\n \\end{matrix} \\right) \\right)\n\\]\non $M_2 \\otimes M_2$\nand is trivial on $C ([0, 1])$.\n(We are only using the coordinates with indexes $1, 2, 3, 4$.)\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is\npurely infinite and simple by Theorem~\\ref{T_0101_PICase},\nbut may well not be ${\\mathcal{Z}}$-stable.\n\\end{exa}\n\nWith a small modification,\nwe can give an example of this type in which the\nunderlying action of $(S^1)^{{\\mathbb{Z}}}$\nis effective.\n\n\\begin{exa}\\label{E_9Z26_OnPINotZStab_v2}\nLet $B$ be the $2^{\\infty}$~UHF algebra,\nand let $\\tau$ be its (unique) tracial state.\nLet ${\\mathrm{tr}}$ be the tracial state{} on~$M_2$,\nand let $\\rho$ be the state on $M_2$ given by\n$\\rho (x)\n = {\\mathrm{tr}} \\big( {\\mathrm{diag}} \\big( \\frac{1}{3}, \\frac{2}{3} \\big) x \\big)$.\nLet $D$ be\nthe reduced free product\n$D = (B \\otimes M_2) \\star_{\\mathrm{r}} C ([0, 1])$,\ntaken with respect to the state $\\tau \\otimes \\rho$ on $B \\otimes M_2$\nand the Lebesgue measure state on $C ([0, 1])$.\nWe claim that $A$ is purely infinite and simple\nbut not ${\\mathcal{Z}}$-stable.\n\nTo prove pure infiniteness,\nin Examples 3.9(iii) of~\\cite{Dkm2}\ntake $A_1 = B$\nwith the state~$\\tau$,\ntake $F = M_2$\nwith the state~$\\rho$,\nand take $B = C ([0, 1])$\nwith the state given by Lebesgue measure.\nThese choices satisfy the hypotheses there.\nSo $A$ is is purely infinite and simple.\n\nTo prove that $A$ is not ${\\mathcal{Z}}$-stable,\nlet $\\pi$ be the Gelfand-Naimark-Segal representation of~$B$\nassociated with~$\\tau$,\nset $N = \\tau (B)''$,\nand also write $\\tau$ for the corresponding tracial state{} on~$N$.\nIn Proposition 5.6 of~\\cite{AmGlJmPh},\ntake $P_1 = N \\otimes M_2$\nwith the state $\\tau \\otimes \\rho$,\nand take $P_2 = L^{\\infty} ([0, 1])$\nwith the state given by Lebesgue measure.\nDefine\n\\[\na = \\left[ \\left( \\begin{matrix} 1 & 0 \\\\ 0 & -1 \\end{matrix} \\right)\n \\otimes 1 \\otimes 1 \\otimes \\cdots \\right]\n \\otimes 1\n \\in \\left( \\bigotimes_{n = 1}^{\\infty} M_2 \\right) \\otimes M_2\n = B \\otimes M_2\n \\S P_1.\n\\]\nTake $G_1 = \\{ 1, a\\} \\S P_1$,\nand take $G_2 \\S P_2$ to be the set of functions\n$\\lambda \\mapsto e^{2 \\pi i n \\lambda}$ for $n \\in {\\mathbb{Z}}$.\nThese choices\nsatisfy the hypotheses there.\nMoreover, $G_1 \\S B \\otimes M_2$ and $G_2 \\S C ([0, 1])$.\nTherefore the reduced free product\n$A = (B \\otimes M_2) \\star_{\\mathrm{r}} C ([0, 1])$\nis a subalgebra of the algebra~$P$\nin Proposition 5.6 of~\\cite{AmGlJmPh}\nwhich contains $a$, $b$, and~$c$,\nso is not ${\\mathcal{Z}}$-stable by Proposition 5.6 of~\\cite{AmGlJmPh}.\nThe claim is proved.\n\nLet $(X, h)$ be\nas in Example~\\ref{E_9214_ShiftAndUHF}\n(and reused in Example~\\ref{E_9214_ShiftAndOI}).\nLet $\\sigma \\colon {\\mathbb{Z}}_{> 0} \\to {\\mathbb{Z}}$ be a bijection\n(as in Example~\\ref{E_9214_ShiftAndUHF}).\nFor $\\zeta = (\\zeta_n)_{n \\in {\\mathbb{Z}}} \\in X$\ntake $\\alpha_{\\zeta}$\nto be the free product automorphism\nwhich is the identity on $C ([0, 1])$\nand which is the infinite tensor product\n\\[\n\\left[ \\bigotimes_{n \\in {\\mathbb{Z}}_{> 0}}\n {\\mathrm{Ad}} \\left( \\left( \\begin{matrix}\n 1 & 0 \\\\\n 0 & \\zeta_{\\sigma (n)}\n \\end{matrix} \\right) \\right) \\right]\n \\otimes {\\mathrm{Ad}} \\left( \\left( \\begin{matrix}\n \\zeta_{\\sigma (1)} & 0 \\\\\n 0 & \\zeta_{\\sigma (2)}\n \\end{matrix} \\right) \\right)\n \\in {\\mathrm{Aut}} (B \\otimes M_2)\n\\]\non~$B \\otimes M_2$.\nThen $C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$ is\npurely infinite and simple by Theorem~\\ref{T_0101_PICase},\nbut may well not be ${\\mathcal{Z}}$-stable.\n\\end{exa}\n\n\\section{Open problems}\\label{Sec_Qs}\n\nIn this section,\nwe collect some open questions\nsuggested by the examples and results in this paper.\n\n\\begin{pbm}\\label{Pb_9922_RokhlinTsr}\nLet $A$ be a unital C*-algebra,\nand let $\\alpha$ be an action of ${\\mathbb{Z}}$ on~$A$\nwhich has finite Rokhlin dimension with commuting towers.\nSuppose that $A$ has stable rank one and $C^* ({\\mathbb{Z}}, A, \\alpha)$ is simple.\nDoes it follow that $C^* ({\\mathbb{Z}}, A, \\alpha)$\nhas stable rank one?\n\\end{pbm}\n\nThis seems to be unknown even if $A$ is simple\n(in which case simplicity of $C^* ({\\mathbb{Z}}, A, \\alpha)$ is automatic)\nand $\\alpha$ has the Rokhlin property.\nWithout assuming simplicity of $C^* ({\\mathbb{Z}}, A, \\alpha)$,\nthe answer is definitely no.\nFor aperiodic homeomorphism{s} of the Cantor set\nwhose transformation group \\ca{s} don't have stable rank one,\nsee Theorem~3.1 of~\\cite{Pt1}\n(it is easy to construct examples there\nwhich have more than one minimal set)\nor Example~8.8 of~\\cite{Phl10}.\nCorollary~2.6 of~\\cite{Szb}\nimplies that the corresponding actions of~${\\mathbb{Z}}$\nhave Rokhlin dimension with commuting towers at most~$1$.\nIn fact, though,\nat least in Example~8.8 of~\\cite{Phl10} we get the Rokhlin property.\n\n\\begin{lem}\\label{L_9922_HasRokhlin}\nLet $h \\colon X \\to X$ be the aperiodic homeomorphism{}\nof the Cantor set in Example~8.8 of~\\cite{Phl10}.\nThen the induced automorphism of $C (X)$ has the Rokhlin property.\n\\end{lem}\n\n\\begin{proof}\nRecall from~\\cite{Phl10} that $X_1 = {\\mathbb{Z}} \\cup \\{\\pm \\infty \\}$,\n$h_1 \\colon X_1 \\to X_1$ is $n \\mapsto n + 1$,\n$X_2$ is the Cantor set,\n$h_2 \\colon X_2 \\to X_2$ is minimal,\n$X = X_1 \\times X_2$,\nand $h = h_1 \\times h_2$.\n\nLet $N \\in {\\mathbb{Z}}_{> 0}$.\nThe standard first return time construction provides\n\\[\nn, r (1), r (2), \\ldots, r (n) \\in {\\mathbb{Z}}_{> 0}\n\\]\nwith $N < r (1) < r (2) < \\cdots < r (n)$,\nand compact open subsets $Z_1, Z_2, \\ldots, Z_n \\S X_2$,\nsuch that\n\\[\nX_2 = \\coprod_{k = 1}^{n} \\coprod_{j = 0}^{r (k) - 1} h_{2}^j (Z_k).\n\\]\nThen the sets\n\\[\nX_1 \\times Z_1, \\quad X_1 \\times Z_2,\n \\quad \\ldots, \\quad X_1 \\times Z_n\n\\]\nform a system of Rokhlin towers for~$h$,\nwith heights $r (1), r (2), \\ldots, r (n)$,\nall of which exceed~$N$.\n\\end{proof}\n\nIn fact,\nit seems to be known that any aperiodic homeomorphism{}\nof the Cantor set~$X$\ninduces an automorphism of $C (X)$ with the Rokhlin property.\nWe have not found a reference,\nand we do not prove this here.\n\nThe following question asks for a plausible generalization\nof Lemma~\\ref{L_4Y11_Comp}.\n\n\\begin{qst}\\label{Qst1_20190103D}\nLet $D$ be a simple unital C*-algebra.\nLet $X$ be a compact metric space, let $G$ be a discrete group,\nand let $(g,x) \\mapsto gx$ be a minimal\nand essentially free action of $G$ on $X$.\nLet $\\alpha \\colon G \\to {\\mathrm{Aut}} (C (X, D))$ be an action of $G$\nwhich lies over the action of $G$ on $X$\nand which is pseudoperiodically generated.\nSet $A = C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$.\n\nDoes it follow that for every $a \\in C (X, D)_{+} \\setminus \\{ 0 \\}$\nthere exists $f \\in C (X)_{+} \\setminus \\{ 0 \\} \\subset A$\nsuch that $f \\precsim_A a$?\n\\end{qst}\n\nWe expect the techniques in the proof of Lemma \\ref{L_4Y11_Comp}\n(taking $Y$ to be the empty set) can be used to answer\nthis question in the affirmative.\n\n\\begin{qst}\\label{Q_9212_ExistsNotPseudoper}\nIs there a simple unital C*-algebra~$D$\nsuch that ${\\mathrm{Aut}} (D)$\nis not pseudoperiodic in the sense of Definition~\\ref{D_4Y12_PsPer}?\n\\end{qst}\n\nPresumably such examples exist,\nbut we don't know of any.\nIndeed,\nwe don't see any reason why there should not be $\\alpha \\in {\\mathrm{Aut}} (D)$\nsuch that $\\set{\\alpha^n \\colon n \\in {\\mathbb{Z}}}$ is not pseudoperiodic.\n\n\\begin{qst}\\label{Q_9214_IsZStab}\nConsider the crossed product\n$C^* \\big( {\\mathbb{Z}}, \\, C (X, E), \\, \\beta \\big)$\nin Example~\\ref{E_9214_ShiftAndCStFI_2}.\nIs this algebra ${\\mathcal{Z}}$-stable?\nWhat about crossed products by similarly constructed actions?\n\\end{qst}\n\n\\begin{qst}\\label{Q_9214_IsZStab_2}\nConsider the crossed product\n$C^* \\big( {\\mathbb{Z}}, \\, C (X)\n \\otimes C^*_{\\mathrm{r}} (F_{\\infty}), \\, \\alpha \\big)$\nin Example~\\ref{E_5416_CSrFGp}.\nIs this algebra ${\\mathcal{Z}}$-stable?\nWhat about crossed products by similarly constructed actions?\nWhat about Example~\\ref{E_5416_DenjFst}?\n\\end{qst}\n\n\\begin{qst}\\label{Q_9217_GenPType_IsZStab}\nConsider the crossed product\n$C^* \\big( {\\mathbb{Z}}, \\, C (S^1, D), \\, \\alpha \\big)$\nin Example~\\ref{E_9217_GenPType}.\nIs this algebra ${\\mathcal{Z}}$-stable?\n\\end{qst}\n\nThe following question\nis motivated by Example~\\ref{E_9214_OnPINotZStab}.\n\n\\begin{qst}\\label{Pb_0111_FinRDim_PI}\nLet $A$ be a nonsimple unital C*-algebra{} which is purely infinite\nin the sense of Definition~4.1 of~\\cite{KiRo2000}.\nLet $\\alpha \\colon {\\mathbb{Z}} \\to {\\mathrm{Aut}} (A)$\nbe an action with finite Rokhlin dimension with commuting towers.\nDoes it follow that $C^* ({\\mathbb{Z}}, A, \\alpha)$\nis purely infinite?\n\\end{qst}\n\nIf $A$ is ${\\mathcal{O}}_{\\infty}$-stable,\nthen $C^* ({\\mathbb{Z}}, A, \\alpha)$ is at least ${\\mathcal{Z}}$-stable\nby Theorem~5.8 of~\\cite{HWZ}.\nIf $A$ is also exact,\nthen so is $C^* ({\\mathbb{Z}}, A, \\alpha)$\n(by Proposition 7.1(v) of~\\cite{Krh2}),\nand $C^* ({\\mathbb{Z}}, A, \\alpha)$ is traceless\n(as at the beginning of Section~5 of~\\cite{Rdm7})\nbecause $A$ is,\nso $C^* ({\\mathbb{Z}}, A, \\alpha)$ is purely infinite\nby Corollary~5.1 of~\\cite{Rdm7}.\n(In this case, one obviously wants ${\\mathcal{O}}_{\\infty}$-stability.\nSee Question~\\ref{Pb_0111_PI_ZStab_to_OIStab} below.)\nBut the question as stated seems to be open,\neven if one assumes that $C^* ({\\mathbb{Z}}, A, \\alpha)$ is simple.\nIf $A$ itself is simple,\nthen $C^* ({\\mathbb{Z}}, A, \\alpha)$\nis purely infinite even just assuming that $\\alpha$ is pointwise outer,\nby Corollary~4.4 of~\\cite{JngOsk}.\nProvided one uses the reduced crossed product,\nthis remains true if ${\\mathbb{Z}}$ is replaced by any discrete group.\n\n\\begin{qst}\\label{Pb_0111_OIStabQ}\nConsider the crossed product\n$C^* \\bigl( {\\mathbb{Z}}, \\, C (X, D), \\, \\alpha \\bigr)$\nin Example~\\ref{E_9214_OnPINotZStab_2}.\nIs this algebra ${\\mathcal{Z}}$-stable?\nIs it ${\\mathcal{O}}_{\\infty}$-stable?\n\\end{qst}\n\nWe hope that ${\\mathcal{Z}}$-stability should come from the\naction in the ``$(S^1)^4$~direction''.\nWe suppose that if a purely infinite simple C*-algebra{}\nis ${\\mathcal{Z}}$-stable,\nthen it is probably ${\\mathcal{O}}_{\\infty}$-stable,\nbut this seems to be open in general.\n\n\\begin{qst}\\label{Pb_0111_PI_ZStab_to_OIStab}\nLet $B$ be a ${\\mathcal{Z}}$-stable purely infinite simple C*-algebra.\nDoes it follow that $B$ is ${\\mathcal{O}}_{\\infty}$-stable?\n\\end{qst}\n\nIf $B$ is separable and nuclear,\nthe answer is yes, by Theorem~5.2 of~\\cite{Rdm7}.\nBut this result doesn't help with Example~\\ref{E_9214_OnPINotZStab_2},\neven if we prove ${\\mathcal{Z}}$-stability there.\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction} \n \\label{sec:Introduction}\n\nThe analysis of gravitational instability in the expanding universe\nmodel was first presented by Lifshitz in 1946 in a general relativistic \ncontext \\cite{Lifshitz}.\nHistorically, the much simpler, and in hindsight, more intuitive \nNewtonian study followed later \\cite{Bonnor}.\nThe pioneering study by Lifshitz is based on a gauge choice which is\ncommonly called the synchronous gauge. \nAs the later studies have shown, the synchronous gauge is only \none way of fixing the gauge freedom out of several available\nfundamental gauge conditions\n\\cite{Harrison,Nariai,Bardeen-1980,KS,Reviews,PRW}.\nAs will be summarized in the following, out of the several gauge conditions\nonly the synchronous gauge fails to fix the gauge mode completely and \nthus often needs more involved algebra.\nAs long as one is careful of the algebra this gauge choice does \nnot cause any kind of intrinsic problem; there exist, however, some persisting\nalgebraic errors in the literature, see Sec. \\ref{sec:Discussion}.\nThe gauge condition which turns out to be especially suitable for\nhandling the perturbed density is the comoving gauge used first by\nNariai in 1969 \\cite{Nariai}.\nSince the comoving gauge condition completely fixes the gauge transformation\nproperty, the variables in this gauge can be equivalently considered as\ngauge-invariant ones.\nAs mentioned, there exist several such fundamental gauge conditions\neach of which completely fixes the gauge transformation properties.\nOne of such gauge conditions suitable for handling the gravitational\npotential and the velocity perturbations is the zero-shear gauge\nused first by Harrison in 1967 \\cite{Harrison}.\nThe variables in such gauge conditions are equivalently gauge-invariant.\nUsing the gauge freedom as an advantage for handling the problem\nwas emphasized by Bardeen in 1988 \\cite{Bardeen-1988}.\nIn order to use the gauge choice as an advantage a {\\it gauge ready method}\nwas proposed in \\cite{PRW} which will be adopted in the \nfollowing, see Sec. \\ref{sec:Relativistic}.\n\nThe variables which characterize the self gravitating Newtonian fluid flow are\nthe density, the velocity and the gravitational potential \n(the pressure is specified by an equation of state),\nwhereas, the relativistic flow may be characterized by various components\nof the metric (and consequent curvatures) and the energy-momentum tensor.\nSince the relativistic gravity theory is a constrained system we have the \nfreedom of imposing certain conditions on the metric or the energy-momentum \ntensor as coordinate conditions. \nIn the perturbation analyses the freedom arises because we need to introduce\na fictitious background system in order to describe the physical perturbed\nsystem.\nThe correspondence of a given spacetime point between the perturbed spacetime\nand the ficticious background one could be made with certain degrees \nof freedom.\nThis freedom can be fixed by suitable conditions (the gauge conditions)\nbased on the spacetime coordinate transformation.\nStudies in the literature show that a certain variable in a certain gauge\ncondition correctly reproduces the corresponding Newtonian behavior.\nAlthough the perturbed density in the comoving gauge shows the Newtonian\nbehavior, the perturbed potential and the perturbed velocity in the\nsame gauge do not behave like the Newtonian ones; for example, in the comoving\ngauge the perturbed velocity vanishes by the coordinate (gauge) condition.\nIt is known that both the perturbed potential and the perturbed velocity\nin the zero-shear gauge correctly behave like the corresponding\nNewtonian variables \\cite{Bardeen-1980}.\n\nIn the first part (Sec. \\ref{sec:Six-gauges}) we will elaborate establishing \nsuch correspondences between the relativistic and Newtonian perturbed \nvariables.\nOur previous work on this subject is presented in \\cite{MDE,Hyun};\ncompared with \\cite{MDE} in this work we will explicitly compare \nthe relativistic equations for the perturbed density, potential and velocity \nvariables in several available gauge conditions with the\nones in the Newtonian system.\nWe will include both the background spatial curvature ($K$) \nand the cosmological constant ($\\Lambda$).\nIn the second part (Sec. \\ref{sec:Hydrodynamics}), using the variables \nwith correct Newtonian correspondences,\nwe will extend our relativistic results to the situations\nwith general pressures in the background and perturbations.\nWe will present the relativistic equations satisfied by the gauge-invariant\nvariables and will derive the {\\it general solutions} valid in the \nsuper-sound-horizon scale (i.e., larger than Jeans scale)\nconsidering both $K$ and $\\Lambda$, and the generally evolving\nideal fluid equation of state $p = p(\\mu)$.\n\nSection \\ref{sec:review} is a review of our previous work\ndisplaying the basic equations in both Newtonian and relativistic contexts\nand summarizing our strategy for handling the equations. \nIn Sec. \\ref{sec:Six-gauges} we consider a pressureless limit of the\nrelativistic equations.\nWe derive the equations for the perturbed density, the perturbed potential \nand the perturbed velocity in several different fundamental gauge conditions.\nBy comparing these relativistic equations in several gauges with\nthe Newtonian ones we {\\it identify} the gauge conditions which reproduce \nthe correct Newtonian behavior for certain variables in general scales.\nIn Sec. \\ref{sec:Hydrodynamics} we present the fully relativistic equations \nfor the identified gauge-invariant variables with correct Newtonian limits, \nnow, considering the general pressures in the background and perturbations.\nWe derive the general large-scale solutions valid for general $K$ and $\\Lambda$\nin an ideal fluid medium, which are thus valid for general equation of \nstate of the \nform $p = p(\\mu)$, but with negligible entropic and anisotropic pressures.\nThe solutions are valid in the super-sound-horizon scale, thus are valid \neffectively in all scales in the matter dominated era where the\nsound-horizon (Jeans scale) is negligible.\nWe discuss several quantities which are conserved in the large-scale\nunder the changing background equation of state.\nThese are the perturbed three-space curvature in several gauge conditions,\nand in particular, we find the three-space curvature perturbation\nin the comoving gauge shows a distinguished behavior:\nfor $K = 0$ (but for general $\\Lambda$), it is conserved in a\nsuper-sound-horizon scale independently of the changing equation of state.\n{}For completeness we also summarize the case with multiple fluid\ncomponents in Sec. \\ref{sec:Multi}, and cases of the rotation and\nthe gravitational wave in Sec. \\ref{sec:rot-GW}.\nSec. \\ref{sec:Hydrodynamics} is the highlight of the present work.\nSec. \\ref{sec:Discussion} is a discussion.\nWe set $c \\equiv 1$.\n\n\n\\section{Basic Equations and Strategy} \n \\label{sec:review}\n\n\n\\subsection{Newtonian Cosmological Perturbations} \n \\label{sec:Newtonian}\n\nSince the Newtonian perturbation analysis in the expanding medium is well known\nwe begin by summarizing the basic equations \\cite{MDE,Weinberg,Peebles-1980}.\nThe background evolution is governed by\n\\begin{eqnarray}\n & & H^2 = {8 \\pi G \\over 3} \\varrho - {K \\over a^2} + {\\Lambda \\over 3}, \n \\quad \\varrho \\propto a^{-3},\n \\label{Newt-BG}\n\\end{eqnarray}\nwhere we allowed the general curvature (total energy) \nand the cosmological constant; $a(t)$ is a cosmic scale factor,\n$H (t) \\equiv \\dot a \/ a$, and $\\varrho (t)$ is the mass density. \nIn Newtonian theory the cosmological constant can be introduced in the\nPoisson equation {\\it by hand} as \n$\\nabla^2 \\Phi = 4 \\pi G \\varrho - \\Lambda$.\nPerturbed parts of the mass conservation, the momentum conservation\nand the Poisson's equations are [see Eqs. (43,46) in \\cite{MDE}]: \n\\begin{eqnarray}\n & & \\delta \\dot \\varrho + 3 H \\delta \\varrho = - {k \\over a} \\varrho\n \\delta v, \\quad\n \\delta \\dot v + H \\delta v = {k \\over a} \\delta \\Phi, \\quad\n - {k^2 \\over a^2} \\delta \\Phi = 4 \\pi G \\delta \\varrho,\n \\label{Newt-eqs} \n\\end{eqnarray}\nwhere $\\delta \\varrho ({\\bf k}, t)$, $\\delta v ({\\bf k}, t)$ and\n$\\delta \\Phi ({\\bf k},t)$ are the Fourier modes of the perturbed mass density, \nvelocity, and gravitational potential, respectively \n[see Eqs. (44,45) in \\cite{MDE}];\n$k$ is a comoving wave number with $\\nabla^2 \\rightarrow - k^2\/a^2$.\n[For linear perturbations the same forms of equations are valid in the \nconfiguration and the Fourier spaces.\nThus, without causing any confusion, we often ignore distinguishing the\nFourier space from the configuration space by an additional subindex.]\nEquation (\\ref{Newt-eqs}) can be arranged into the closed form equations\nfor $\\delta$ ($\\equiv \\delta \\varrho \/ \\varrho$), $\\delta v$ and\n$\\delta \\Phi$ as:\n\\begin{eqnarray}\n & & \\ddot \\delta + 2 H \\dot \\delta - 4 \\pi G \\varrho \\delta \n = {1 \\over a^2 H} \\left[ a^2 H^2 \\left( {\\delta \\over H} \\right)^\\cdot\n \\right]^\\cdot = 0,\n \\label{Newt-delta-eq} \\\\\n & & \\delta \\ddot v + 3 H \\delta \\dot v + \\left( \\dot H + 2 H^2 \n - 4 \\pi G \\varrho \\right) \\delta v = 0,\n \\label{Newt-v-eq} \\\\\n & & \\delta \\ddot \\Phi + 4 H \\delta \\dot \\Phi\n + \\left( \\dot H + 3 H^2 - 4 \\pi G \\varrho \\right) \\delta \\Phi \n = {1 \\over a^3 H} \\left[ a^2 H^2 \n \\left( {a \\over H} \\delta \\Phi \\right)^\\cdot \\right]^\\cdot = 0.\n \\label{Newt-Phi-eq} \n\\end{eqnarray}\nWe note that Eqs. (\\ref{Newt-eqs}-\\ref{Newt-Phi-eq}) are\nvalid for general $K$ and $\\Lambda$.\nAlthough redundant, we presented these equations for later comparison\nwith the relativistic results.\nThe general solutions for $\\delta$, $\\delta v$ and $\\delta \\Phi$\nimmediately follow as (see also Table 1 in \\cite{MDE}):\n\\begin{eqnarray}\n & & \\delta ({\\bf k},t) = k^2 \\left[ \n H C ({\\bf k}) \\int^t_0 {dt \\over \\dot a^2} \n + {H \\over 4 \\pi G \\varrho a^3} d ({\\bf k}) \\right], \n \\label{delta-N} \\\\\n & & \\delta v ({\\bf k}, t) = - \\left[ {k \\over aH} C ({\\bf k})\n \\left( 1 + a^2 H \\dot H \\int^t_0 {dt \\over \\dot a^2} \\right)\n + {k \\dot H \\over 4 \\pi G \\varrho a^2} d ({\\bf k}) \\right],\n \\label{v-N} \\\\\n & & \\delta \\Phi ({\\bf k}, t) = - \\left[ 4 \\pi G \\varrho a^2 H C ({\\bf k}) \n \\int^t_0 {dt \\over \\dot a^2} + {H \\over a} d ({\\bf k}) \\right].\n \\label{Phi-N}\n\\end{eqnarray}\nThe coefficients $C ({\\bf k})$ and $d ({\\bf k})$ are two integration constants, \nand indicate the relatively growing and decaying solutions, respectively;\nthe coefficients are matched in accordance with the solutions with a general\npressure in Eqs. (\\ref{varphi_chi-sol}-\\ref{v_chi-sol}) using \nEq. (\\ref{Correspondences-2})\n(different dummy variables for $t$ in the integrands are assumed).\nIn the literature we can find various analytic form expressions of the \nabove solutions in diverse situations, \nsee \\cite{Lifshitz,Harrison,Nariai,Weinberg,Peebles-1980,MDE-sols}.\n\n\n\\subsection{Relativistic Cosmological Perturbations} \n \\label{sec:Relativistic}\n\nIn the relativistic theory the fundamental variables are the metric and\nthe energy-momentum tensor.\nAs a way of summarizing our notation we present our convention of\nthe metric and the energy-momentum tensor.\nAs the metric we consider a spatially homogeneous and isotropic spacetime\nwith the most general perturbation\n\\begin{eqnarray}\n d s^2 = - a^2 \\left( 1 + 2 \\alpha \\right) d \\eta^2\n - a^2 \\left( \\beta_{,\\alpha} +B_\\alpha \\right) d \\eta d x^\\alpha\n + a^2 \\left[ g_{\\alpha\\beta}^{(3)} \\left( 1 + 2 \\varphi \\right)\n + 2 \\gamma_{|\\alpha\\beta} + 2 C_{(\\alpha|\\beta)} + 2 C_{\\alpha\\beta}\n \\right] d x^\\alpha d x^\\beta,\n \\label{metric} \n\\end{eqnarray}\nwhere $0 = \\eta$ with $dt \\equiv a d \\eta$,\nIndices $\\alpha$, $\\beta$ $\\dots$ are spatial ones and\n$g^{(3)}_{\\alpha\\beta}$ is the three-space metric of the homogeneous and\nisotropic space.\n$\\alpha ({\\bf x}, t)$, $\\beta ({\\bf x}, t)$, $\\gamma ({\\bf x}, t)$,\nand $\\varphi ({\\bf x}, t)$ indicate the scalar-type structure with\nfour degrees of freedom.\nThe transverse $B_\\alpha ({\\bf x}, t)$ and $C_\\alpha ({\\bf x}, t)$\nindicate the vector-type structure with four degrees of freedom.\nThe transverse-tracefree $C_{\\alpha\\beta} ({\\bf x}, t)$\nindicates the tensor-type structure with two degrees of freedom.\nThe energy-momentum tensor is decomposed as:\n\\begin{eqnarray}\n & & T^0_0 \\equiv - \\bar \\mu - \\delta \\mu, \\quad\n T^0_\\alpha \\equiv - {1 \\over k} \n \\left( \\bar \\mu + \\bar p \\right) v_{,\\alpha}\n + \\delta T^{(v)0}_{\\;\\;\\;\\;\\;\\alpha},\n \\nonumber \\\\\n & & T^\\alpha_\\beta\n \\equiv \\left( \\bar p + \\delta p \\right) \\delta^\\alpha_\\beta \n + {1 \\over a^2} \\left( \\nabla^\\alpha \\nabla_\\beta\n - {1 \\over 3} \\Delta \\delta^\\alpha_\\beta \\right) \\sigma\n + \\delta T^{(v)\\alpha}_{\\;\\;\\;\\;\\;\\beta}\n + \\delta T^{(t)\\alpha}_{\\;\\;\\;\\;\\;\\beta},\n \\label{Tab}\n\\end{eqnarray}\nwhere $\\mu ({\\bf k}, t) \\equiv \\bar \\mu (t) + \\delta \\mu ({\\bf k}, t)$ and \n$p \\equiv \\bar p + \\delta p$ are the energy density and the pressure, \nrespectively; an overbar indicates a background order quantity and will \nbe ignored unless necessary.\n$v$ is a frame-independent velocity related perturbed order variable\nand $\\sigma$ is an anisotropic pressure.\nSituation with multiple fluids will be considered in Sec. \\ref{sec:Multi}.\n$\\delta T^{(v)0}_{\\;\\;\\;\\;\\;\\alpha}$, $\\delta T^{(v)\\alpha}_{\\;\\;\\;\\;\\;\\beta}$,\nand $\\delta T^{(t)\\alpha}_{\\;\\;\\;\\;\\;\\beta}$ are vector and tensor-type\nperturbed order energy-momentum tensor.\nAll the spatial tensor indices and the vertical bar are based on \n$g^{(3)}_{\\alpha\\beta}$ as the metric.\nThe three types of structures are related to the density condensation,\nthe rotation, and the gravitational wave, respectively.\nDue to the high symmetry in the background the three types of perturbations\ndecouple from each other to the linear order.\nSince these three types of structures evolve independently, we can handle\nthem separately.\nEvolutions of the rotation and the gravitational wave will be discussed\nin Sec. \\ref{sec:rot-GW}.\n\n{}For the background equations we have [Eq. (6) in \\cite{Ideal}]:\n\\begin{eqnarray}\n H^2 = {8 \\pi G \\over 3} \\mu - {K \\over a^2} + {\\Lambda \\over 3}, \\quad\n \\dot H = - 4 \\pi G ( \\mu + p ) + {K \\over a^2}, \\quad\n \\dot \\mu = - 3 H \\left( \\mu + p \\right). \n \\label{BG-eqs} \n\\end{eqnarray}\nThese equations follow from $G^0_0$ and $G^\\alpha_\\alpha - 3 G^0_0$\ncomponents of Einstein equations and $T^b_{0;b} = 0$, respectively;\nthe third equation follows from the first two equations.\n{}For $p = 0$ and replacing $\\mu$ with $\\varrho$ Eq. (\\ref{BG-eqs}) \nreduces to Eq. (\\ref{Newt-BG}).\nThe relativistic cosmological perturbation equations without fixing\nthe temporal gauge condition, thus in a gauge ready form,\nwere derived in Eq. (22-28) of \\cite{PRW}\n[see also Eqs. (8-14) in \\cite{Ideal}; we use \n$\\delta \\equiv \\delta \\mu \/\\mu$ and $v \\equiv - (k\/a) \\Psi\/(\\mu + p)$]:\n\\begin{eqnarray}\n & & \\kappa \\equiv 3 \\left( - \\dot \\varphi + H \\alpha \\right)\n + {k^2 \\over a^2} \\chi,\n \\label{eq1} \\\\\n & & - {k^2 - 3K \\over a^2} \\varphi + H \\kappa = - 4 \\pi G \\mu \\delta,\n \\label{eq2} \\\\\n & & \\kappa - {k^2 - 3K \\over a^2} \\chi = 12 \\pi G \\left( \\mu + p \\right)\n {a \\over k} v,\n \\label{eq3} \\\\\n & & \\dot \\chi + H \\chi - \\alpha - \\varphi = 8 \\pi G \\sigma,\n \\label{eq4} \\\\\n & & \\dot \\kappa + 2 H \\kappa = \\left( {k^2 \\over a^2} - 3 \\dot H \\right) \n \\alpha + 4 \\pi G \\left( 1 + 3 c_s^2 \\right) \\mu \\delta \n + 12 \\pi G e,\n \\label{eq5} \\\\\n & & \\dot \\delta + 3 H \\left( c_s^2 - {\\rm w} \\right) \\delta \n + 3 H { e \\over \\mu } = \\left( 1 + {\\rm w} \\right)\n \\left( \\kappa - 3 H \\alpha - {k \\over a} v \\right),\n \\label{eq6} \\\\\n & & \\dot v + \\left( 1 - 3 c_s^2 \\right) H v = {k \\over a} \\alpha\n + {k \\over a \\left( 1 + {\\rm w} \\right) } \\left( c_s^2 \\delta\n + {e \\over \\mu} - {2\\over 3} {k^2 - 3K \\over a^2} {\\sigma \\over \\mu}\n \\right),\n \\label{eq7} \n\\end{eqnarray}\nwhere we introduced a spatially gauge-invariant combination\n$\\chi \\equiv a ( \\beta + a \\dot \\gamma)$.\nAccording to their origins, Eqs. (\\ref{eq1}-\\ref{eq7}) can be called: \nthe definition of $\\kappa$, \nADM energy constraint ($G^0_0$ component of Einstein equations), \nmomentum constraint ($G^0_\\alpha$ component), \nADM propagation \n($G^\\alpha_\\beta - {1 \\over 3} \\delta^\\alpha_\\beta G^\\gamma_\\gamma$ component),\nRaychaudhuri ($G^\\alpha_\\alpha - G^0_0$ component), \nenergy conservation ($T^b_{0;b} = 0$),\nand momentum conservation ($T^b_{\\alpha;b} = 0$) equations, respectively.\nThe perturbed metric variables $\\varphi ({\\bf k}, t)$, $\\kappa ({\\bf k},t)$,\n$\\chi ({\\bf k}, t)$ and $\\alpha ({\\bf k}, t)$ are the perturbed part of\nthe three space curvature, expansion, shear, and the lapse function, \nrespectively; meanings for $\\varphi$, $\\kappa$, and $\\chi$ are based on\nthe normal frame vector field.\nThe isotropic pressure is decomposed as\n\\begin{eqnarray}\n \\delta p ({\\bf k}, t) \\equiv c_s^2 (t) \\delta \\mu ({\\bf k}, t) \n + e ({\\bf k}, t), \\quad\n c_s^2 \\equiv {\\dot p \\over \\dot \\mu}, \\quad\n {\\rm w} (t) \\equiv {p \\over \\mu}.\n \\label{p-decomposition}\n\\end{eqnarray}\n\nUnder the gauge transformation $\\tilde x^a = x^a + \\xi^a$, the variables\ntransform as (see Sec. 2.2 in \\cite{PRW}):\n\\begin{eqnarray}\n & & \\tilde \\alpha = \\alpha - \\dot \\xi^t, \\quad\n \\tilde \\varphi = \\varphi - H \\xi^t, \\quad\n \\tilde \\chi = \\chi - \\xi^t, \\quad\n \\tilde \\kappa = \\kappa \n + \\left( 3 \\dot H - {k^2 \\over a^2} \\right) \\xi^t,\n \\nonumber \\\\\n & & \\tilde v = v - {k \\over a} \\xi^t, \\quad\n \\tilde \\delta = \\delta + 3 \\left( 1 + {\\rm w} \\right) H \\xi^t,\n \\label{GT}\n\\end{eqnarray}\nwhere $\\xi^t \\equiv e^s \\xi^\\eta$, and $e$ and $\\sigma$ are \ngauge-invariant.\n[Due to the spatial homogeneity of the background, the effect \n from the spatial gauge transformation has been trivially handled; \n $\\chi$ and $v$ are the spatially gauge-invariant combinations of \n the variables, and the other metric and fluid variables \n are naturally spatially gauge-invariant \\cite{Bardeen-1988}.]\nThus, the perturbed variables in Eqs. (\\ref{eq1}-\\ref{eq7})\nare {\\it designed} so that any one of the following conditions\ncan be used to fix the freedom based on the temporal gauge transformation:\n$\\alpha \\equiv 0$ (the synchronous gauge), \n$\\varphi \\equiv 0$ (the uniform-curvature gauge),\n$\\chi \\equiv 0$ (the zero-shear gauge), \n$\\kappa \\equiv 0$ (the uniform-expansion gauge), \n$v\/k \\equiv 0$ (the comoving gauge), \nand $\\delta \\equiv 0$ (the uniform-density gauge).\nThese gauge conditions contain most of the gauge conditions used in the\nliterature concerning the cosmological hydrodynamic perturbation in a single\ncomponent situation; for the multiple component case see Eq. (\\ref{Gauges-2}).\nOf course, we can make an infinite number of different linear combinations \nof these gauge conditions each of which is also suitable for the temporal \ngauge fixing condition.\nOur reason for choosing these conditions as fundamental is partly based \non conventional use, but apparently, as the names of the conditions indicate,\nthe gauge conditions also have reasonable meanings based on the metric and \nthe energy-momentum tensors.\n\nEach one of these gauge conditions, except for the synchronous gauge, \nfixes the temporal gauge transformation property completely.\nThus, each variable in these five gauge conditions \nis equivalent to a corresponding gauge-invariant combination\n(of the variable concerned and the variable used in the gauge condition).\nWe proposed a convenient way of writing the gauge-invariant variables \n\\cite{PRW}.\n{}For example, we let\n\\begin{eqnarray}\n & & \\delta_v \\equiv \\delta + 3 (1 + {\\rm w}) {aH \\over k} v\n \\equiv 3 (1 + {\\rm w}) {a H \\over k} v_\\delta, \\quad\n \\varphi_\\chi \\equiv \\varphi - H \\chi \\equiv - H \\chi_\\varphi, \\quad\n v_\\chi \\equiv v - {k \\over a} \\chi \\equiv - {k \\over a} \\chi_v.\n \\label{GI-combinations}\n\\end{eqnarray}\nThe variables $\\delta_v$, $v_\\chi$ and $\\varphi_\\chi$ are gauge-invariant \ncombinations; $\\delta_v$ becomes $\\delta$ in the comoving gauge which takes\n$v\/k=0$ as the gauge condition, etc.\nIn this manner we can systematically construct the corresponding \ngauge-invariant combination for any variable based on a gauge condition which \nfixes the temporal gauge transformation property completely.\nOne variable evaluated in different gauges can be considered as different \nvariables, and they show different behaviors in general.\nIn this sense, except for the synchronous gauge, the variables in the \nrest of five gauges can be considered as the gauge-invariant variables.\n[The variables with a subindex $\\alpha$ are not gauge-invariant, \nbecause those are equivalent to variables in the synchronous gauge.]\nAlthough $\\delta_v$ is a gauge-invariant variable which becomes\n$\\delta$ in the comoving gauge, $\\delta_v$ itself \ncan be evaluated in any gauge with the same value.\n\nEquations (\\ref{eq1}-\\ref{eq7}) are the basic set of gauge-ready form\nequations for hydrodynamic perturbations proposed in \\cite{PRW}: \n``The moral is that one should work in the gauge that is mathematically\nmost convenient for the problem at hand'', \\cite{Bardeen-1988}.\nAccordingly, in order to use the available gauge conditions as the\nadvantage, Eqs. (\\ref{eq1}-\\ref{eq7}) are designed \nfor easy implementation of the fundamental gauge conditions.\nUsing this {\\it gauge-ready formulation}, complete sets of solutions for the \nsix different gauge conditions are presented in an ideal fluid case \n\\cite{Ideal}, and in a pressureless medium \\cite{MDE}.\nEquations (\\ref{BG-eqs}-\\ref{p-decomposition}) include general pressures which \nmay account for the nonequilibrium or dissipative effects in hydrodynamic \nflows in the cosmological context with general $K$ and $\\Lambda$.\nThe equations are expressed in general forms so that the fluid quantities \ncan represent successfully the effects of the scalar field and \nclasses of generalized gravity theories \\cite{PRW}.\nThe applications of the gauge-ready formalism to the minimally coupled scalar\nfield and to the generalized gravity theories have been made in \\cite{MSF-GGT}.\n\n\n\\section{Newtonian Correspondence of the Relativistic Analyses} \n \\label{sec:Six-gauges}\n\nIn \\cite{MDE,Hyun} we developed arguments on the correspondences\nbetween the Newtonian and the relativistic analyses.\nIn order to reinforce the Newtonian correspondences of certain\n(gauge-invariant) variables in certain gauges, in the following we\nwill present the closed form differential equations for \n$\\delta$, $v$ and $\\varphi$ in the six different gauge conditions\nand compare them with the Newtonian ones in \nEqs. (\\ref{Newt-delta-eq}-\\ref{Newt-Phi-eq}).\n\nThere are reasons why we should know about possible Newtonian behaviors\nof the relativistic variables.\nVariables in the relativitistic gravity are only parameters appearing in\nthe metric and energy-momentum tensor.\n{}For example, later we will find that only in the zero-shear gauge does\nthe perturbed curvature variable behave as the perturbed Newtonian potential,\nwhich we have some experience of in Newtonian physics.\nThe same variable evaluated in other gauge conditions is simply other\nvariable and we cannot regard it as the perturbed potential.\nAs we have introduced several different fundamental gauge conditions\nit became necessary to identify the correct gauge where a variable shows\nthe corresponding Newtonian behavior.\nAs will be shown in this section, we rarely find Newtonian correspondences.\nIn fact, there exists an almost unique gauge condition for each\nrelativistic perturbation variable which shows Newtonian behavior.\nThe results will be summarized in Sec. \\ref{sec:Correspondences}.\n\nIn this section we consider a {\\it pressureless medium},\n$p = 0$ (thus ${\\rm w} = 0 = c_s^2$) and $e = 0 = \\sigma$.\nHowever, we consider general $K$ and $\\Lambda$.\n\n\n\\subsection{Comoving Gauge} \n \\label{sec:CG}\n\nAs the gauge condition we set $v\/k \\equiv 0$.\nEquivalently, we can set $v\/k = 0$ and leave the other variables as\nthe gauge-invariant combinations with subindices $v\/k$ or simply $v$.\n{}From Eq. (\\ref{eq7}) we have $\\alpha_v = 0$.\nThus the comoving gauge is a case of the synchronous gauge; this is\ntrue {\\it only} in a pressureless situation.\n{}From Eqs. (\\ref{eq1}-\\ref{eq7}) we can derive:\n\\begin{eqnarray}\n & & \\ddot \\delta_v + 2 H \\dot \\delta_v - 4 \\pi G \\mu \\delta_v = 0,\n \\label{CG-delta-eq} \\\\\n & & \\ddot \\varphi_v + 3 H \\dot \\varphi_v - {K \\over a^2} \\varphi_v = 0.\n \\label{CG-varphi-eq} \n\\end{eqnarray}\nThus, Eq. (\\ref{CG-delta-eq}) has a form identical to Eq. \n(\\ref{Newt-delta-eq}).\nHowever, Eq. (\\ref{CG-varphi-eq}) differs from Eq. (\\ref{Newt-Phi-eq}),\nand apparently we do not have an equation for $v$.\n\n{}For $K = 0$ Eq. (\\ref{CG-varphi-eq}) leads to\ntwo solutions which are $\\varphi_v \\propto {\\rm constant}$ and \n$\\int_0^t a^{-3} dt$. \nSince we have $\\alpha_v = 0$, for $K = 0$, from Eqs. (\\ref{eq1},\\ref{eq3})\nwe have $\\dot \\varphi_v = 0$.\nThis implies that, for $K = 0$, the second solution, \n$\\varphi_v \\propto \\int_0^t a^{-3} dt$, should have the vanishing coefficient,\nsee Eq. (\\ref{varphi_v-sol3}).\nThe solution of Eq. (\\ref{CG-varphi-eq}) for general $K$ will be\npresented later in Eq. (\\ref{varphi_v-sol}).\n\n\n\\subsection{Synchronous Gauge} \n \\label{sec:SG}\n\nWe let $\\alpha \\equiv 0$.\nThis gauge condition does not fix the temporal gauge transformation property\ncompletely.\nThus, although we can still indicate the variables in this gauge condition\nusing subindices $\\alpha$ without ambiguity, these variables are not \ngauge-invariant; see the discussion after Eq. (\\ref{varphi-GI}).\nEquation (\\ref{eq7}) leads to\n\\begin{eqnarray}\n v_\\alpha = c_g {k \\over a},\n\\end{eqnarray}\nwhich is a pure gauge mode.\nThus, fixing $c_g \\equiv 0$ exactly corresponds to the comoving gauge.\nWe can show that the following two equations are not affected by the\nremaining gauge mode.\n\\begin{eqnarray}\n & & \\ddot \\delta_\\alpha + 2 H \\dot \\delta_\\alpha \n - 4 \\pi G \\mu \\delta_\\alpha = 0,\n \\label{SG-delta-eq} \\\\\n & & \\ddot \\varphi_\\alpha + 3 H \\dot \\varphi_\\alpha \n - {K \\over a^2} \\varphi_\\alpha = 0.\n \\label{SG-varphi-eq} \n\\end{eqnarray}\nEquation (\\ref{SG-delta-eq}) is identical to Eq. (\\ref{Newt-delta-eq}).\nThis is because in a pressureless medium the behavior of the gauge modes \nhappen to coincide with the behavior of the decaying physical solutions\nfor $\\delta_\\alpha$ and $\\varphi_\\alpha$.\nHowever, for the variables $\\kappa_\\alpha$ and $\\chi_\\alpha$ the gauge mode \ncontribution appears explicitly.\n\nThus, in a pressureless medium, variables\nin the synchronous gauge behave the same as the ones in the comoving gauge, \nexcept for the additional gauge modes which appear in the synchronous gauge\nfor some variables.\nIn a pressureless medium, we can simultaneously impose both the comoving gauge\nand the synchronous gauge conditions.\nHowever, this is possible only in a pressureless medium, see the Appendix.\nIn Sec. \\ref{sec:Discussion} we indicate some common errors in the literature\nbased on the synchronous gauge.\n\n\n\\subsection{Zero-shear Gauge} \n \\label{sec:ZSG}\n\nWe let $\\chi \\equiv 0$, and substitute the other variables into the \ngauge-invariant combinations with subindices $\\chi$.\nWe can derive:\n\\begin{eqnarray}\n & & \\ddot \\delta_\\chi + { 2 k^2 \/ a^2 - 36 \\pi G \\mu \\left[ 1 + ( H^2 \n - 2 \\dot H ) a^2 \/ \\left( k^2 - 3 K \\right) \\right] \\over\n k^2 \/ a^2 - 12 \\pi G \\mu \\left[ 1 + 3 H^2 a^2\/\\left( k^2 - 3K \\right) \n \\right] }\n H \\left( \\dot \\delta_\\chi \n - 12 \\pi G \\mu {a^2 \\over k^2 - 3 K} H \\delta_\\chi \\right)\n \\nonumber \\\\\n & & \\qquad \n - 4 \\pi G \\mu \\left[ 1 - 3 ( H^2 - 2 \\dot H )\n {a^2 \\over k^2 - 3 K} \\right] \\delta_\\chi = 0,\n \\label{ZSG-delta-eq} \\\\\n & & \\ddot v_\\chi + 3 H \\dot v_\\chi \n + \\left( \\dot H + 2 H^2 - 4 \\pi G \\mu \\right) v_\\chi = 0,\n \\label{ZSG-v-eq} \\\\\n & & \\ddot \\varphi_\\chi + 4 H \\dot \\varphi_\\chi + \\left( \\dot H + 3 H^2 \n - 4 \\pi G \\mu \\right) \\varphi_\\chi = 0.\n \\label{ZSG-varphi-eq} \n\\end{eqnarray}\nThus, Eqs. (\\ref{ZSG-v-eq},\\ref{ZSG-varphi-eq}) have forms identical \nto Eqs. (\\ref{Newt-v-eq},\\ref{Newt-Phi-eq}), respectively.\nHowever, only in the small-scale limit is the behavior of $\\delta_\\chi$\nthe same as the Newtonian one.\n\n\n\\subsection{Uniform-expansion Gauge} \n \\label{sec:UEG}\n\nWe let $\\kappa \\equiv 0$, and substitute the remaining variables into\nthe gauge-invariant combinations with subindices $\\kappa$.\nWe can derive:\n\\begin{eqnarray}\n & & \\ddot \\delta_\\kappa \n + \\left( 2 H - {12 \\pi G \\mu H \\over k^2\/a^2 - 3 \\dot H}\n \\right) \\dot \\delta_\\kappa - \\left[ 4 \\pi G \\mu {k^2\/a^2 + 6 H^2 \\over\n k^2\/a^2 - 3 \\dot H} + \\left( {12 \\pi G \\mu H \\over\n k^2\/a^2 - 3 \\dot H} \\right)^\\cdot \\right] \\delta_\\kappa = 0,\n \\label{UEG-delta-eq} \\\\ \n & & \\ddot v_\\kappa + \\left[ 2 H - {12 \\pi G \\mu H \\over k^2\/a^2 - 3 \\dot H}\n + {4 \\pi G \\mu \\over k^2\/a^2 - 3 \\dot H}\n \\left( {k^2\/a^2 - 3 \\dot H \\over 4 \\pi G \\mu} \\right)^\\cdot\n \\right] \\dot v_\\kappa\n \\nonumber \\\\ \n & & \\quad \n + \\left[ \\dot H + H^2 - 4 \\pi G \\mu {k^2\/a^2 + 3 H^2 \\over\n k^2\/a^2 - 3 \\dot H} + H {4 \\pi G \\mu \\over k^2\/a^2 - 3 \\dot H}\n \\left( { k^2\/a^2 - 3 \\dot H \\over 4 \\pi G \\mu} \\right)^\\cdot\n \\right] v_\\kappa = 0,\n \\label{UEG-v-eq} \\\\ \n & & \\ddot \\varphi_\\kappa \n + \\left( 4 H - {12 \\pi G \\mu H \\over k^2\/a^2 - 3 \\dot H}\n \\right) \\dot \\varphi_\\kappa\n + \\left[ \\dot H + 3 H^2 - 4 \\pi G \\mu {k^2\/a^2 + 9 H^2 \\over\n k^2\/a^2 - 3 \\dot H} - \\left( {12 \\pi G \\mu H \\over k^2\/a^2\n - 3 \\dot H} \\right)^\\cdot \\right] \\varphi_\\kappa = 0.\n \\label{UEG-varphi-eq}\n\\end{eqnarray}\nIn the small-scale limit we can show that\nEqs. (\\ref{UEG-delta-eq}-\\ref{UEG-varphi-eq}) \nreduce to Eqs. (\\ref{Newt-delta-eq}-\\ref{Newt-Phi-eq}), respectively.\nThus, {\\it in the small-scale limit}, all three variables $\\delta_\\kappa$, \n$v_\\kappa$ and $\\varphi_\\kappa$ correctly reproduce the Newtonian behavior.\nHowever, outside or near the (visual) horizon scale, the behaviors of all these\nvariables {\\it strongly deviate} from the Newtonian ones.\n\nIn Sec. 84 of \\cite{Peebles-1980} we find that in order to get the usual \nNewtonian equations a coordinate transformation is made so that\nwe have $\\dot h \\equiv 0$ in the new coordinate.\nWe can show that $\\dot h = 2 \\kappa$ in our notation.\nThus, the new coordinate used in \\cite{Peebles-1980} is the uniform-expansion \ngauge\\footnote{This was incorrectly pointed out after Eq. (49) in \\cite{MDE}:\n In \\cite{MDE} it was wishfully mentioned that in Sec. 84 of \n \\cite{Peebles-1980} the gauge transformations were made into the \n comoving gauge for $\\delta$ and into the zero-shear gauge for \n $v$ and $\\varphi$, respectively.\n }.\n\n\n\\subsection{Uniform-curvature Gauge} \n \\label{sec:UCG}\n\nWe let $\\varphi \\equiv 0$, and substitute the other variables into the \ngauge-invariant combinations with subindices $\\varphi$.\nWe have\n\\begin{eqnarray}\n & & \\ddot \\delta_\\varphi + 2 H { (k^2 - 3K)\/a^2 + 18 \\pi G \\mu \\over \n (k^2 - 3K)\/a^2 + 12 \\pi G \\mu } \\dot \\delta_\\varphi\n - { 4 \\pi G \\mu k^2\/a^2 \\over \n (k^2 - 3K)\/a^2 + 12 \\pi G \\mu } \\delta_\\varphi = 0,\n \\label{UCG-delta-eq} \\\\\n & & \\ddot v_\\varphi + \\left( 5 H + 2 {\\dot H \\over H} \\right) \\dot v_\\varphi\n + \\left( 3 \\dot H + 4 H^2 \\right) v_\\varphi = 0. \n \\label{UCG-v-eq}\n\\end{eqnarray}\nIn the small-scale limit Eq. (\\ref{UCG-delta-eq}) reduces to\nEq. (\\ref{Newt-delta-eq}).\nIn the uniform-curvature gauge, the perturbed potential is set to be equal\nto zero by the gauge condition.\nThe uniform-curvature gauge condition is known to have \ndistinct properties in handling the scalar field or the dilaton field \nwhich appears in a broad class of the generalized gravity theories\n\\cite{MSF-GGT,GRG-MSF-GGT}.\n\n\n\\subsection{Uniform-density Gauge} \n \\label{sec:UDG}\n\nWe let $\\delta \\equiv 0$, and substitute the other variables into the \ngauge-invariant combinations with subindices $\\delta$.\nWe have\n\\begin{eqnarray}\n & & \\ddot v_\\delta + 2 \\left( 2 H + {\\dot H \\over H} \\right) \\dot v_\\delta\n + 3 \\left( H^2 + \\dot H \\right) v_\\delta = 0,\n \\label{UDG-v-eq} \\\\\n & & \\ddot \\varphi_\\delta + 2 H { (k^2 - 3 K)\/a^2 + 18 \\pi G \\mu\n \\over (k^2- 3K)\/a^2 + 12 \\pi G \\mu} \\dot \\varphi_\\delta\n - {4 \\pi G \\mu k^2\/a^2\n \\over (k^2- 3K)\/a^2 + 12 \\pi G \\mu} \\varphi_\\delta = 0.\n \\label{UDG-varphi-eq} \n\\end{eqnarray}\nThese equations differ from Eqs. (\\ref{Newt-v-eq},\\ref{Newt-Phi-eq}).\nOf course, we have no equation for $\\delta$ which is set to be equal to zero \nby our choice of the gauge condition.\n\n\n\\subsection{Newtonian Correspondences} \n \\label{sec:Correspondences}\n\nAfter a thorough comparison made in this section we found that \nequations for $\\delta$ in the comoving gauge ($\\delta_v$), and\nfor $v$ and $\\varphi$ in the zero-shear gauge ($v_\\chi$ and $\\varphi_\\chi$)\nshow the same forms as the corresponding Newtonian equations \n{\\it in general scales} and considering general $K$ and $\\Lambda$.\nUsing these gauge-invariant combinations in Eq. (\\ref{GI-combinations}),\nEqs. (\\ref{eq1}-\\ref{eq7}) can be combined to give:\n\\begin{eqnarray}\n \\dot \\delta_v = - {k^2 - 3K \\over a k} v_\\chi, \\quad\n \\dot v_\\chi + H v_\\chi = - {k \\over a} \\varphi_\\chi, \\quad\n {k^2 - 3K \\over a^2} \\varphi_\\chi = 4 \\pi G \\mu \\delta_v.\n \\label{GI-eqs}\n\\end{eqnarray}\nComparing Eq. (\\ref{GI-eqs}) with Eq. (\\ref{Newt-eqs}) \nwe can identify either one of the following correspondences:\n\\begin{eqnarray}\n & & \\delta_v \\leftrightarrow \\delta |_{\\rm Newtonian}, \\quad\n {k^2 - 3 K \\over k^2} v_\\chi \n \\leftrightarrow \\delta v |_{\\rm Newtonian}, \\quad\n - {k^2 - 3 K \\over k^2} \\varphi_\\chi \\leftrightarrow \n \\delta \\Phi |_{\\rm Newtonian},\n \\label{Correspondences-1} \\\\\n & & {k^2 \\over k^2 - 3 K} \\delta_v \\leftrightarrow \n \\delta |_{\\rm Newtonian}, \\quad\n v_\\chi \\leftrightarrow \\delta v |_{\\rm Newtonian}, \\quad\n - \\varphi_\\chi \\leftrightarrow \\delta \\Phi |_{\\rm Newtonian}.\n \\label{Correspondences-2}\n\\end{eqnarray}\nIn \\cite{MDE} we proposed the correspondences in Eq. (\\ref{Correspondences-1}),\nbut the ones in Eq. (\\ref{Correspondences-2}) also work well.\nIn fact, the gravitational potential identified in \nEq. (\\ref{Correspondences-2}) is the one often found in the Newtonian limit \nof the general relativity, e.g. in \\cite{Weinberg}.\nUsing Eqs. (\\ref{GI-combinations},\\ref{eq3},\\ref{eq4},\\ref{GI-eqs}) \nwe can also identify the following relations:\n\\begin{eqnarray}\n & & \\delta_v = 3 H {a \\over k} v_\\delta, \\quad\n v_\\chi = - {k \\over a} \\chi_v = - {ak \\over k^2 - 3 K} \\kappa_v, \\quad\n \\varphi_\\chi = -\\alpha_\\chi = - H \\chi_\\varphi.\n\\end{eqnarray}\n\n{}From a given solution known in one gauge we can derive all the rest of the \nsolutions even in other gauge conditions.\nThis can be done either by using the gauge-invariant combination of variables \nor directly through gauge transformations.\nThe complete set of exact solutions in a pressureless medium is presented \nin \\cite{MDE}.\n{}From our study in this section and using the complete solutions presented \nin Tables 1 and 2 of \\cite{MDE} we can identify variables in certain gauges \nwhich correspond to the Newtonian ones. These are summarized in Table 1.\n\n\\centerline{ [[{\\bf TABLE 1.}]] }\n\n\\noindent\nAs mentioned earlier, all three variables in the uniform-expansion gauge show \nNewtonian correspondence in the small-scale; however all of them change \nthe behaviors near and above the horizon scale. \nIn the small-scale limit, except for the \nuniform-density gauge where $\\delta \\equiv 0$, \n$\\delta$ in all gauge conditions behaves in the same way \\cite{Bardeen-1980}.\nMeanwhile, since only $\\delta_v$ ($\\delta$ in the comoving gauge)\nshows the Newtonian behavior in a general scale, it may be natural to identify \n$\\delta_v$ as the one most closely resembling the Newtonian one.\n\nNotice that, although we have horizons in the relativistic analysis\nthe equations for $\\delta_v$, $v_\\chi$ and $\\varphi_\\chi$\nkeep the same form as the corresponding Newtonian equations\n{\\it in the general scale}.\nConsidering this as an additional point we regard\n$\\delta_v$, $v_\\chi$ and $\\varphi_\\chi$ as the one most closely corresponding\nto the Newtonian variables.\nThis argument will become stronger as we consider the case with a general\npressure in the next section.\n\n\n\\section{Relativistic Cosmological Hydrodynamics}\n \\label{sec:Hydrodynamics}\n\n\n\\subsection{General Equations and Large-scale Solutions}\n \\label{sec:Eqs}\n\nIn the previous sections we have shown that the gauge-invariant\ncombinations $\\delta_v$, $v_\\chi$ and $\\varphi_\\chi$ behave most similarly\nto the Newtonian $\\delta \\equiv \\delta \\varrho\/\\varrho$, \n$\\delta v$ and $\\delta \\Phi$, respectively.\nThe equations remain the same in a general scale which includes the\nsuperhorizon scales in the relativistic situation considering \ngeneral $K$ and $\\Lambda$.\nIn this section, we will present the fully general relativistic equations \nfor $\\delta_v$, $v_\\chi$ and $\\varphi_\\chi$ including the effects of the \ngeneral pressure terms.\nEquations (\\ref{eq1}-\\ref{eq7}) are the basic set of \nperturbation equations in a gauge-ready form.\n\n{}From Eqs. (\\ref{eq3},\\ref{eq6},\\ref{eq7}),\nEqs. (\\ref{eq4},\\ref{eq7}), Eqs. (\\ref{eq2},\\ref{eq3}), and\nEqs. (\\ref{eq1},\\ref{eq3},\\ref{eq4}), we have, respectively:\n\\begin{eqnarray}\n & & \\dot \\delta_v - 3 H {\\rm w} \\delta_v\n = - \\left( 1 + {\\rm w} \\right) {k \\over a} {k^2 - 3 K \\over k^2} v_\\chi\n - 2 H {k^2 - 3K \\over a^2} {\\sigma \\over \\mu},\n \\label{P-1} \\\\\n & & \\dot v_\\chi + H v_\\chi = - {k \\over a} \\varphi_\\chi\n + {k \\over a \\left( 1 + {\\rm w} \\right)} \n \\left[ c_s^2 \\delta_v + {e \\over \\mu}\n - 8 \\pi G \\left( 1 + {\\rm w} \\right) \\sigma\n - {2 \\over 3} {k^2 - 3 K \\over a^2} {\\sigma \\over \\mu} \\right],\n \\label{P-2} \\\\\n & & {k^2 - 3 K \\over a^2} \\varphi_\\chi = 4 \\pi G \\mu \\delta_v,\n \\label{P-3} \\\\\n & & \\dot \\varphi_\\chi + H \\varphi_\\chi\n = - 4 \\pi G \\left( \\mu + p \\right) {a \\over k} v_\\chi \n - 8 \\pi G H \\sigma.\n \\label{P-4}\n\\end{eqnarray}\nConsidering either one of the correspondences in \nEqs. (\\ref{Correspondences-1},\\ref{Correspondences-2})\nwe immediately see that Eqs. (\\ref{P-1}-\\ref{P-4}) have the \ncorrect Newtonian limit expressed in Eqs. (\\ref{Newt-eqs},\\ref{GI-eqs}).\nEquations (\\ref{P-1}-\\ref{P-4}) were presented in \\cite{Bardeen-1980}.\n\nCombining Eqs. (\\ref{P-1}-\\ref{P-4}) we can derive closed\nform expressions for the $\\delta_v$ and $\\varphi_\\chi$\nwhich are the relativistic counterpart of Eqs. \n(\\ref{Newt-delta-eq},\\ref{Newt-Phi-eq}).\nWe have:\n\\begin{eqnarray}\n & & \\ddot \\delta_v + \\left( 2 + 3 c_s^2 - 6 {\\rm w} \\right) H \\dot \\delta_v\n + \\Bigg[ c_s^2 {k^2 \\over a^2} \n - 4 \\pi G \\mu \\left( 1 - 6 c_s^2 + 8 {\\rm w} - 3 {\\rm w}^2 \\right)\n + 12 \\left( {\\rm w} - c_s^2 \\right) {K \\over a^2}\n + \\left( 3 c_s^2 - 5 {\\rm w} \\right) \\Lambda \\Bigg] \\delta_v\n \\nonumber \\\\\n & & \\qquad \\qquad\n = {1 + {\\rm w} \\over a^2 H} \\left[ {H^2 \\over a (\\mu + p)} \n \\left( {a^3 \\mu \\over H} \\delta_v \\right)^\\cdot \\right]^\\cdot\n + c_s^2 {k^2 \\over a^2} \\delta_v\n \\nonumber \\\\\n & & \\qquad \\qquad\n = - {k^2 - 3 K \\over a^2} { 1 \\over \\mu }\n \\left\\{ e + 2 H \\dot \\sigma\n + 2 \\left[ - {1 \\over 3} {k^2 \\over a^2} + 2 \\dot H\n + 3 \\left( 1 + c_s^2 \\right) H^2 \\right] \\sigma \\right\\},\n \\label{delta-eq}\\\\\n & & \\ddot \\varphi_\\chi + \\left( 4 + 3 c_s^2 \\right) H \\dot \\varphi_\\chi\n + \\left[ c_s^2 {k^2 \\over a^2} \n + 8 \\pi G \\mu \\left( c_s^2 - {\\rm w} \\right)\n - 2 \\left( 1 + 3 c_s^2 \\right) {K \\over a^2}\n + \\left( 1 + c_s^2 \\right) \\Lambda \\right] \\varphi_\\chi\n \\nonumber \\\\\n & & \\qquad \\qquad\n = {\\mu + p \\over H} \\left[ {H^2 \\over a (\\mu + p)} \n \\left( {a \\over H} \\varphi_\\chi \\right)^\\cdot \\right]^\\cdot\n + c_s^2 {k^2 \\over a^2} \\varphi_\\chi\n = {\\rm stresses}.\n \\label{varphi-eq}\n\\end{eqnarray}\nThese two equations are related by Eq. (\\ref{P-3}) which resembles the \nPoisson equation.\nNotice the remarkably compact forms presented in the second steps of\nEqs. (\\ref{delta-eq},\\ref{varphi-eq}).\nIt may be an interesting exercise to show that the above equations\nare indeed valid for general $K$ and $\\Lambda$, and\nfor the general equation of state $p = p(\\mu)$;\nuse $\\dot {\\rm w} = - 3 H (1 + {\\rm w}) ( c_s^2 - {\\rm w} )$.\nEquation (\\ref{delta-eq}) became widely known through Bardeen's seminal paper \nin \\cite{Bardeen-1980}.\nIn a less general context but originally \nEqs. (\\ref{delta-eq},\\ref{varphi-eq}) were derived by\nNariai \\cite{Nariai} and Harrison \\cite{Harrison}, respectively;\nhowever, the compact expressions are new results in this paper.\n\nIf we ignore the entropic and anisotropic pressures ($e = 0 = \\sigma$) \non scales larger than Jeans scale\nEq. (\\ref{varphi-eq}) immediately leads to a general integral form \nsolution as\n\\begin{eqnarray}\n \\varphi_\\chi ({\\bf k}, t)\n &=& 4 \\pi G C ({\\bf k}) {H \\over a} \\int_0^t {a (\\mu + p) \\over H^2} dt\n + {H \\over a} d ({\\bf k})\n \\nonumber \\\\\n &=& C ({\\bf k}) \\left[ 1 - {H \\over a} \\int_0^t\n a \\left( 1 - {K \\over \\dot a^2} \\right) dt \\right]\n + {H \\over a} d ({\\bf k}).\n \\label{varphi_chi-sol}\n\\end{eqnarray}\nThis solution was first derived in Eq. (108) of \n\\cite{HV}\\footnote{We correct a typographical \n error in Eq. (129) of \\cite{HV}: $c_0\/a$ should be replaced by $c_0$.\n },\nsee also Eq. (55) in \\cite{PRW}.\nSolutions for $\\delta_v$ and $v_\\chi$ follow from Eqs.\n(\\ref{P-3},\\ref{P-4}), respectively, as\n\\begin{eqnarray}\n & & \\delta_v ({\\bf k}, t) = {k^2 - 3 K \\over 4 \\pi G \\mu a^2} \n \\varphi_\\chi ({\\bf k}, t), \n \\label{delta_v-sol} \\\\\n & & v_\\chi ({\\bf k}, t)= - {k \\over 4 \\pi G ( \\mu + p) a^2}\n \\left\\{ C ({\\bf k}) \\left[ {K \\over \\dot a} \n - \\dot H \\int_0^t a \\left( 1 - {K \\over \\dot a^2} \\right) dt \\right]\n + \\dot H d ({\\bf k}) \\right\\}.\n \\label{v_chi-sol}\n\\end{eqnarray}\nWe stress that these large-scale asymptotic solutions are\n{\\it valid for the general $K$, $\\Lambda$, and $p = p(\\mu)$};\n$C({\\bf k})$ and $d({\\bf k})$ are integration constants considering \nthe general evolution of $p = p(\\mu)$.\nWe also emphasize that {\\it the large-scale criterion is the sound-horizon},\nand thus in the matter dominated era the solutions are valid even far \ninside the (visual) horizon as long as the scales are larger than Jeans scale.\nIn the pressureless limit, considering Eq. (\\ref{Correspondences-2}),\nEqs. (\\ref{varphi_chi-sol}-\\ref{v_chi-sol})\ncorrectly reproduce Eqs. (\\ref{delta-N}-\\ref{Phi-N}).\n\nUsing the set of gauge-ready form equations in\nEqs. (\\ref{eq1}-\\ref{eq7}) the complete set of corresponding general\nsolutions for all variables in all different gauges can be easily derived;\nfor such sets of solutions with less general assumptions see \\cite{MDE,Ideal}.\n{}For $K = 0 = \\Lambda$ and ${\\rm w} = {\\rm constant}$ we have\n$a \\propto t^{2\/[3(1+{\\rm w})]}$ and Eqs. \n(\\ref{varphi_chi-sol}-\\ref{v_chi-sol}) become:\n\\begin{eqnarray}\n \\varphi_\\chi ({\\bf k}, t) \n &=& {3 (1 + {\\rm w}) \\over 5 + 3 {\\rm w}} C ({\\bf k})\n + {2 \\over 3 (1 + {\\rm w})} {1 \\over at} d ({\\bf k})\n \\nonumber \\\\\n &\\propto& {\\rm constant}, \n \\quad t^{- {5 + 3 {\\rm w} \\over 3 ( 1 + {\\rm w} )} }\n \\quad \\propto \\quad {\\rm constant}, \n \\quad a^{- {5 + 3 {\\rm w} \\over 2} },\n \\nonumber \\\\\n \\delta_v ({\\bf k}, t) \n &\\propto& t^{2(1 + 3 {\\rm w}) \\over 3 (1 + {\\rm w})}, \n \\quad t^{- {1 - {\\rm w} \\over 1 + {\\rm w} } }\n \\quad \\propto \\quad a^{1 + 3 {\\rm w}}, \n \\quad a^{- {3 \\over 2} (1 - {\\rm w}) },\n \\nonumber \\\\\n v_\\chi ({\\bf k}, t)\n &\\propto& t^{1 + 3 {\\rm w} \\over 3 (1 + {\\rm w})}, \n \\quad t^{- {4 \\over 3(1 + {\\rm w}) } }\n \\quad \\propto \\quad a^{1 + 3 {\\rm w} \\over 2}, \n \\quad a^{- 2}.\n \\label{Sol-w=const}\n\\end{eqnarray}\nThese solutions were presented in \\cite{Bardeen-1980} and\nin less general contexts but originally in \\cite{Harrison,Nariai}.\n{}For ${\\rm w} =0$ we recover the well known Newtonian behaviors.\n \n\n\\subsection{Conserved Quantities}\n \\label{sec:Conservation}\n\nIn an ideal fluid, the curvature fluctuations in several gauge conditions \nare known to be conserved in the {\\it superhorizon scale} independently\nof the changes in the background equation of state.\n{}From Eqs. (41,73) in \\cite{Ideal} we find [for $K = 0 = \\Lambda$,\nbut for general $p(\\mu)$]:\n\\begin{eqnarray}\n \\varphi_v = \\varphi_\\delta = \\varphi_\\kappa = \\varphi_\\alpha = C ({\\bf k}),\n \\label{varphi-conserv}\n\\end{eqnarray}\nand the dominating decaying solutions vanish (or, are cancelled).\nThe combinations follow from Eq. (\\ref{GT}) as:\n\\begin{eqnarray}\n \\varphi_v = \\varphi - {aH \\over k} v, \\quad\n \\varphi_\\delta \\equiv \\varphi +{\\delta \\over 3 (1 + {\\rm w})} \n \\equiv \\zeta, \\quad\n \\varphi_\\kappa \\equiv \\varphi + {H \\over 3 \\dot H - k^2\/a^2 } \\kappa, \n \\quad\n \\varphi_\\alpha \\equiv \\varphi - H \\int^t \\alpha dt.\n \\label{varphi-GI}\n\\end{eqnarray}\nThe $\\varphi_\\alpha$ combination which is $\\varphi$ in the synchronous gauge\nis not gauge-invariant; the lower bound of integration gives the behavior\nof the gauge mode; in Eq. (\\ref{varphi-conserv}) we ignored the gauge mode.\nThe large-scale conserved quantity, $\\zeta$, proposed in\n\\cite{Bardeen-1988,BST} is $\\varphi_\\delta$.\nIn \\cite{Ideal} the above conservation properties are shown assuming \n$K = 0 = \\Lambda$; in such a case the integral form solutions for a complete\nset of variables are presented in Table 8 in \\cite{Ideal}.\nIn Eqs. (\\ref{varphi_chi-sol}-\\ref{v_chi-sol}) we have the large-scale \nintegral form solutions in the case of general $K$ and $\\Lambda$.\nThus, now we can see the behavior of these variables considering the \ngeneral $K$ and $\\Lambda$.\n\n{}Evaluating $\\varphi_v$ in Eq. (\\ref{varphi-GI}) in the zero-shear gauge,\nthus $\\varphi_v = \\varphi_\\chi - (aH \/ k) v_\\chi$, and using\nthe solutions in Eqs. (\\ref{varphi_chi-sol},\\ref{v_chi-sol}) we \ncan derive\n\\begin{eqnarray}\n \\varphi_v ({\\bf k}, t) \n &=& C ({\\bf k}) \\left\\{ 1 + {K \\over a^2}\n {1 \\over 4 \\pi G (\\mu + p)} \\left[ 1 \n - {H \\over a} \\int_0^t a \\left( 1 - {K \\over \\dot a^2} \\right) dt \n \\right] \\right\\}\n + {K \\over a^2} { H\/a \\over 4 \\pi G ( \\mu + p)} d ({\\bf k})\n \\nonumber \\\\\n &=& C ({\\bf k}) \\left[ 1 + {K \\over a^2}\n {H\/a \\over \\mu + p} \\int_0^t {a (\\mu + p) \\over H^2} dt \\right]\n + {K \\over a^2} { H\/a \\over 4 \\pi G ( \\mu + p)} d ({\\bf k}).\n \\label{varphi_v-sol}\n\\end{eqnarray}\nThus, for $K = 0$ (but for general $\\Lambda$) we have\n\\begin{eqnarray}\n \\varphi_v ({\\bf k}, t) = C ({\\bf k}),\n \\label{varphi_v-sol2}\n\\end{eqnarray}\nwith the vanishing decaying solution. \nThe disappearance of the decaying \nsolution in Eq. (\\ref{varphi_v-sol2}) implies that the dominating decaying \nsolutions in Eqs. (\\ref{varphi_chi-sol},\\ref{v_chi-sol}) cancel out \nfor $K = 0$.\nIn fact, for $K = 0$, from Eqs. \n(\\ref{eq1},\\ref{eq2},\\ref{eq3},\\ref{eq6},\\ref{eq7}) we can derive\n\\begin{eqnarray}\n & & {c_s^2 H^2 \\over a^3 ( \\mu + p)} \\left[ {a^3 ( \\mu + p) \\over c_s^2 H^2}\n \\dot \\varphi_v \\right]^\\cdot\n + c_s^2 {k^2 \\over a^2} \\varphi_v = {\\rm stresses}.\n \\label{varphi_v-fluid-eq}\n\\end{eqnarray}\n{}For a pressureless case ($c_s^2 = 0$), instead of\nEq. (\\ref{varphi_v-fluid-eq}) we have $\\dot \\varphi_v = 0$;\nsee after Eq. (\\ref{CG-varphi-eq}).\nIn the large-scale limit, and ignoring the stresses, we have an integral\nform solution\n\\begin{eqnarray}\n & & \\varphi_v\n = C ({\\bf x}) - \\tilde D ({\\bf x})\n \\int^t_0 {c_s^2 H^2 \\over a^3 ( \\mu + p ) } dt.\n \\label{varphi_v-sol3}\n\\end{eqnarray}\nCompared with the solution in Eq. (\\ref{varphi_v-sol2}), the $\\tilde D$ \nterm in Eq. (\\ref{varphi_v-sol3}) is $c_s^2(k\/aH)^2$ order higher\nthan $d$ terms in the other gauge.\nTherefore, for $K = 0$, $\\varphi_v$ is conserved for the generally varying\nbackground equation of state, i.e., for general $p = p(\\mu)$.\nSince the solutions in \nEqs. (\\ref{varphi_chi-sol},\\ref{delta_v-sol},\\ref{v_chi-sol}) are valid \nin super-sound-horizon scale, the conservation property in \nEq. (\\ref{varphi_v-sol2}) is valid in all scales in the matter \ndominated era.\n\n{}For $\\varphi_\\delta$, by evaluating it in the comoving gauge we have\n\\begin{eqnarray}\n \\varphi_\\delta \\equiv \\varphi_v + {\\delta_v \\over 3 (1 + {\\rm w})}.\n\\end{eqnarray}\nUsing Eqs. (\\ref{varphi_v-sol},\\ref{delta_v-sol}) we have\n\\begin{eqnarray}\n \\varphi_\\delta ({\\bf k}, t)\n &=& C ({\\bf k}) \\left\\{ 1 + {k^2 \\over a^2}\n {1 \\over 12 \\pi G (\\mu + p)} \\left[ 1\n - {H \\over a} \\int_0^t a \\left( 1 - {K \\over \\dot a^2} \\right) dt\n \\right] \\right\\}\n + {k^2 \\over a^2} { H\/a \\over 12 \\pi G ( \\mu + p)} d ({\\bf k}).\n \\label{varphi_delta-sol}\n\\end{eqnarray}\nSince the higher order term ignored in Eq. (\\ref{varphi_v-sol}) \nis $c_s^2 (k\/aH)^2$ order higher, in the medium with $c_s \\sim 1 (\\equiv c)$,\nthe terms involving $(k\/a)^2\/(\\mu + p)$ are not necessarily valid.\nThus, to the leading order in the superhorizon scale we have\n\\begin{eqnarray}\n \\varphi_\\delta ({\\bf k}, t) = C ({\\bf k}).\n\\end{eqnarray}\nThus, in the superhorizon scale $\\varphi_\\delta$ is conserved \napparently considering $K$ and $\\Lambda$.\n[In the case of nonvanishing $K$ we may need to be \n careful in defining the superhorizon scale. \n Here, as the superhorizon condition we simply {\\it took} $k^2$ term \n to be negligible.\n In a hyperbolic (negatively curved) space with $K = -1$, \n since $R^{(3)} = {6 K \/ a^2}$, the distance \n $a\/\\sqrt{|K|}$ introduces a curvature scale: on smaller scales the \n space is near flat and on larger scales the curvature term dominates.\n In the unit of the curvature scale $k^2$ less than one ($\\sim |K|$)\n corresponds to the super-curvature scale, and $k^2$ larger than \n one corresponds to the sub-curvature scale. \n Thus, when we ignored the $k^2$ term compared with $|K|$ in \n a negatively curved space, we are considering the super-curvature scale.\n In other words, for the hyperbolic backgrounds $k^2$ takes continuous \n value with $k^2 \\ge 0$ where $k^2 \\ge 1$ corresponds to subcurvature \n mode whereas $0\\le k^2 <1 $ corresponds to supercurvature mode, and \n our large-scale limit took $k^2 \\rightarrow 0$.\n A useful study in the hyperbolic situation can be found \n in \\cite{Lyth-etal}.]\nHowever, in a medium with the negligible sound speed (like the matter dominated \nera), Eq. (\\ref{varphi_delta-sol}) is valid considering the $k^2$ terms.\nThus, as the scale enters the horizon in the matter dominated era \n$\\varphi_\\delta$ wildly changes its behavior and is {\\it no longer conserved \ninside the horizon};\ninside the horizon $\\varphi_\\delta$ is dominated by the $\\delta$ part\nwhich shows Newtonian behavior of $\\delta$ in most of the gauge conditions, \nsee Table 1.\n\n{}For $\\varphi_\\kappa$, by evaluating it in the uniform-density gauge\nand using Eq. (\\ref{eq2}), which gives \n$\\kappa_\\delta = (k^2 - 3K)\/(a^2 H) \\varphi_\\delta$, we have\n\\begin{eqnarray}\n \\varphi_\\kappa = \\varphi_\\delta\n \\left[ 1 + {k^2 - 3K \\over 12 \\pi G (\\mu + p) a^2} \\right]^{-1}.\n \\label{varphi_kappa-sol}\n\\end{eqnarray}\nThus, in the superhorizon scale and for $K = 0$, $\\varphi_\\kappa$ is conserved.\nAs in $\\varphi_\\delta$, in a medium with the negligible sound speed, \nthe solution in Eq. (\\ref{varphi_kappa-sol}) with Eq. (\\ref{varphi_delta-sol}) \nis valid considering the $k^2$ terms.\nThus, as the scale enters the horizon in the matter dominated era \n$\\varphi_\\kappa$\nalso changes its behavior and is {\\it no longer conserved inside the horizon}.\n\nNow we summarize:\nIn the superhorizon scale $\\varphi_v$ and $\\varphi_\\kappa$ are conserved for \n$K = 0$, while $\\varphi_\\delta$ is conserved considering the general $K$.\nIn the matter dominated era, for $K = 0$, $\\varphi_v$ is still conserved in the\nsuper-sound-horizon scale which virtually covers all scales,\nwhereas $\\varphi_\\delta$ and $\\varphi_\\kappa$ change their behaviors\nnear the horizon crossing epoch and are no longer conserved inside the horizon.\nIn this regard, we may say, $\\varphi_v$ is the {\\it best conserved quantity}.\nCuriously, although $\\varphi_\\chi$ is the variable most closely resembling\nthe Newtonian gravitational potential, it is not conserved for changing \nequation of state, see Eq. (\\ref{varphi_chi-sol});\nhowever, it is conserved independently of the horizon crossing in the\nmatter dominated era for $K = 0 = \\Lambda$, see Eq. (\\ref{Sol-w=const}).\n\nSimilar conservation properties of the curvature variable in certain gauge\nconditions remain true for models based on a minimally coupled scalar \nfield and even on classes of generalized gravity theories.\nIn the generalized gravity the uniform-field gauge is more suitable\nfor handling the conservation property, and the uniform-field gauge\ncoincides with the comoving gauge in the minimally coupled scalar field.\n{}Thorough studies of the minimally coupled scalar field\nand the generalized gravity, assuming $K = 0 = \\Lambda$, are made in \n\\cite{MSF-GGT,MSF-GGT-2} and summarized in \\cite{GRG-MSF-GGT}.\n\n\n\\subsection{Multi-component situation}\n \\label{sec:Multi}\n\nWe consider the energy-momentum tensor composed of multiple components as\n\\begin{eqnarray}\n T_{ab} = \\sum_{(l)} T_{(l)ab}, \\quad\n T^{\\;\\;\\;\\; b}_{(i)a;b} \\equiv Q_{(i)a}, \\quad\n \\sum_{(l)} Q_{(l)a} = 0,\n \\label{Tab-multi}\n\\end{eqnarray}\nwhere $i,l (= 1,2, \\dots, n)$ indicates $n$ fluid components, and\n$Q_{(i)a}$ considers possible mutual interactions among fluids.\nSince we are considering the scalar-type perturbation we decompose \nthe $Q_{(i)a}$ in the following way:\n\\begin{eqnarray}\n Q_{(i)0} \\equiv - a (1 + \\alpha) Q_{(i)}, \\quad\n Q_{(i)} \\equiv \\bar Q_{(i)} + \\delta Q_{(i)}, \\quad\n Q_{(i)\\alpha} \\equiv J_{(i),\\alpha}.\n\\end{eqnarray}\nThe scalar-type fluid quantities in Eq. (\\ref{Tab}) can be considered \nas the collective ones which are related to the individual component as:\n\\begin{eqnarray}\n & & \\bar \\mu = \\sum_{(l)} \\bar \\mu_{(l)}, \\quad\n \\bar p = \\sum_{(l)} \\bar p_{(l)}; \\quad\n \\delta \\mu = \\sum_{(l)} \\delta \\mu_{(l)}, \\quad\n \\delta p = \\sum_{(l)} \\delta p_{(l)}, \n \\nonumber \\\\\n & & (\\mu + p) v = \\sum_{(l)} (\\mu_{(l)} + p_{(l)} ) v_{(l)}, \\quad\n \\sigma = \\sum_{(l)} \\sigma_{(l)}.\n \\label{fluid-multi}\n\\end{eqnarray}\n\n{}For the background Eq. (\\ref{BG-eqs}) remains valid for the collective fluid\nquantities.\nAn additional equation follows from the individual\nenergy-momentum conservation in Eq. (\\ref{Tab-multi}) as\n\\begin{eqnarray}\n \\dot \\mu_{(i)} + 3 H \\left( \\mu_{(i)} + p_{(i)} \\right) = Q_{(i)}.\n \\label{BG-multi}\n\\end{eqnarray}\n{}For the perturbations Eqs. (\\ref{eq1}-\\ref{eq7}) remain valid for the \ncollective fluid quantities. \nThe additional equations we need follow from the individual\nenergy-momentum conservation in Eq. (\\ref{Tab-multi}).\n{}From $T^{\\;\\;\\;\\; b}_{(i)0;b} = Q_{(i)0}$ and using Eq. (\\ref{eq1})\nwe have the energy conservation of the fluid components \n\\begin{eqnarray}\n \\delta \\dot \\mu_{(i)} + 3 H \\left( \\delta \\mu_{(i)}\n + \\delta p_{(i)} \\right)\n = - {k \\over a} \\left( \\mu_{(i)} + p_{(i)} \\right) v_{(i)}\n + \\dot \\mu_{(i)} \\alpha\n + \\left( \\mu_{(i)} + p_{(i)} \\right) \\kappa + \\delta Q_{(i)}.\n \\label{eq8}\n\\end{eqnarray}\n{}From $T^{\\;\\;\\;\\; b}_{(i)\\alpha;b} = Q_{(i)\\alpha}$ we have\nthe momentum conservation of the fluid components\n\\begin{eqnarray}\n \\dot v_{(i)} + \\left[ H \\left( 1 - 3 c_{(i)}^2 \\right)\n + { 1 + c_{(i)}^2 \\over \\mu_{(i)} + p_{(i)} } Q_{(i)} \\right] v_{(i)}\n = {k \\over a} \\left[ \\alpha\n + {1 \\over \\mu_{(i)} + p_{(i)}} \\left( \\delta p_{(i)}\n - {2 \\over 3} {k^2 - 3 K \\over a^2} \\sigma_{(i)}\n - J_{(i)} \\right) \\right].\n \\label{eq9}\n\\end{eqnarray}\nBy adding Eqs. (\\ref{eq8},\\ref{eq9}) over the components we have\nEqs. (\\ref{eq6},\\ref{eq7}).\nIn the multi-component situation \nEqs. (\\ref{eq1}-\\ref{eq7},\\ref{eq8},\\ref{eq9}) are the complete set \nexpressed in a gauge-ready form.\nIf we decompose the pressure as\n$\\delta p_{(i)} = c_{(i)}^2 \\delta \\mu_{(i)} + e_{(i)}$ with\n$c_{(i)}^2 \\equiv \\dot p_{(i)}\/\\dot \\mu_{(i)}$, comparing with\nEq. (\\ref{p-decomposition}) we have\n\\begin{eqnarray}\n e = \\delta p - c_s^2 \\delta \\mu\n = \\sum_{(l)} \\left[ \\left( c_{(l)}^2 - c_s^2 \\right) \\delta \\mu_{(l)}\n + e_{(l)} \\right].\n \\label{e}\n\\end{eqnarray}\n\nUnder the gauge transformation, from Eq. (20) in \\cite{PRW}, we have:\n\\begin{eqnarray}\n & & \\delta \\tilde \\mu_{(i)} = \\delta \\mu_{(i)} - \\dot \\mu_{(i)} \\xi^t, \\quad\n \\delta \\tilde p_{(i)} = \\delta p_{(i)} - \\dot p_{(i)} \\xi^t, \\quad\n \\tilde v_{(i)} = v_{(i)} - {k \\over a} \\xi^t, \n \\nonumber \\\\\n & & \\tilde J_{(i)} = J_{(i)} + Q_{(i)} \\xi^t, \\quad\n \\delta \\tilde Q_{(i)} = \\delta Q_{(i)} - \\dot Q_{(i)} \\xi^t.\n \\label{GT-2}\n\\end{eqnarray}\nThus, we have the following additional gauge conditions:\n\\begin{eqnarray}\n \\delta \\mu_{(i)} \\equiv 0, \\quad\n \\delta p_{(i)} \\equiv 0, \\quad\n v_{(i)}\/k \\equiv 0, \\quad\n {\\rm etc}.\n \\label{Gauges-2}\n\\end{eqnarray}\nAny one of these gauge conditions also fixes the temporal gauge condition\ncompletely.\nStudies of the multi-component situations\ncan be found in \\cite{KS,PRW,KS-2}.\n\n\n\\subsection{Rotation and Gravitational Wave}\n \\label{sec:rot-GW}\n\n{}For the vector-type perturbation, using Eqs. (\\ref{metric},\\ref{Tab}),\nwe have:\n\\begin{eqnarray}\n & & 8 \\pi G T^{(v)0}_{\\;\\;\\;\\;\\;\\alpha}\n = - {\\Delta + 2 K \\over 2 a^2} \\left( B_\\alpha + a \\dot C_\\alpha \\right),\n \\label{rot-1} \\\\\n & & 8 \\pi G \\delta T^{(v)\\alpha}_{\\;\\;\\;\\;\\;\\beta}\n = {1 \\over 2 a^3} \\left\\{ a^2 \\left[ B^\\alpha_{\\;\\;|\\beta}\n + B_\\beta^{\\;\\;|\\alpha} + a \\left( C^\\alpha_{\\;\\;|\\beta}\n + C_\\beta^{\\;\\;|\\alpha} \\right)^\\cdot \\right] \\right\\}^\\cdot,\n \\label{rot-2} \\\\\n & & {1 \\over a^3} \\left( a^4 T^{(v)0}_{\\;\\;\\;\\;\\;\\alpha} \\right)^\\cdot\n = - \\delta T^{(v)\\beta}_{\\;\\;\\;\\;\\;\\alpha|\\beta},\n \\label{rot-3}\n\\end{eqnarray}\nwhere Eq. (\\ref{rot-3}) follows from $T^b_{\\alpha ; b} = 0$, and \napparently Eqs. (\\ref{rot-1},\\ref{rot-2}) follow from the Einstein equations.\nWe can show that Eq. (\\ref{rot-3}) follows from Eqs. (\\ref{rot-1},\\ref{rot-2}).\nWe introduce a gauge-invariant velocity related variable \n$V^{(\\omega)}_\\alpha ({\\bf x}, t)$ as \n$T^{(v)0}_{\\;\\;\\;\\;\\;\\alpha} \\equiv (\\mu + p) V^{(\\omega)}_\\alpha$.\nWe can show $V^{(\\omega)}_\\alpha \\propto a \\omega$ where $\\omega$ \nis the amplitude of vorticity vector, see Sec. 5 of \\cite{PRW}.\nThus, ignoring the anisotropic stress in the RHS of Eq. (\\ref{rot-3})\nwe have the conservation of angular momentum as\n\\begin{eqnarray}\n {\\rm Angular \\; Momentum} \\sim a^3 \\left( \\mu + p \\right)\n \\times a \\times V^{(\\omega)}_\\alpha ({\\bf x}, t)\n = {\\rm constant \\; in \\; time},\n\\end{eqnarray}\nwhich is valid for general $K$, $\\Lambda$, and $p(\\mu)$ in general scale.\n\n{}For the tensor-type perturbation, using Eqs. (\\ref{metric},\\ref{Tab})\nin the Einstein equations, we have\n\\begin{eqnarray}\n 8 \\pi G \\delta T^{(t)\\alpha}_{\\;\\;\\;\\;\\;\\beta}\n = \\ddot C^\\alpha_\\beta + 3 H \\dot C^\\alpha_\\beta\n - {\\Delta - 2 K \\over a^2} C^\\alpha_\\beta.\n \\label{GW-eq}\n\\end{eqnarray}\nIn the large-scale limit, assuming $K = 0$ and ignoring the anisotropic\npressure Eq. (\\ref{GW-eq}) has a general integral form solution\n\\begin{eqnarray}\n C^\\alpha_\\beta ({\\bf x}, t)\n = C^{\\;\\alpha}_{1\\beta} ({\\bf x}) - D^\\alpha_\\beta ({\\bf x})\n \\int^t {1 \\over a^3} dt,\n\\end{eqnarray}\nwhere $C^{\\;\\alpha}_{1\\beta}$ and $D^\\alpha_\\beta$ are integration constants\nfor relatively growing and decaying solutions, respectively.\nThus, ignoring the transient term the amplitude of the gravitational\nwave is conserved in the super-horizon scale considering general\nevolution of the background equation of state.\n\nThe conservation of angular momentum and the equation for the\ngravitational wave were first derived in \\cite{Lifshitz}.\nEvolutions in the multi-component situation were considered in\nSec. 5 of \\cite{PRW}.\n\n\n\\section{Discussions}\n \\label{sec:Discussion}\n\nWe would like to make comments on related works in several {\\it textbooks}:\nEq. (15.10.57) in \\cite{Weinberg}, Eq. (10.118) in \\cite{Peebles-2},\nEq. (11.5.2) in \\cite{Coles}, Eq. (8.52) in \\cite{Moss},\nand problem 6.10 in \\cite{Padmanabhan} are in error.\nAll these errors are essentially about the same point involved with\na fallible algebraic mistake in the synchronous gauge.\nThe correction in the case of \\cite{Weinberg} was made in \\cite{GRG}:\nin a medium with a nonvanishing pressure the equation for the density \nfluctuation in the synchronous gauge becomes third order because of \nthe presence of a gauge mode in addition to the physical growing \nand decaying solutions which is true even in the large-scale limit.\nThe truncated second order equation in \\cite{Weinberg,Moss} \npicks up a gauge mode instead of the physical decaying solution in the \nsynchronous gauge.\nThe errors in \\cite{Peebles-2,Padmanabhan} are based on imposing \nthe synchronous gauge and the comoving gauge simultaneously, and thus \nhappen to end up with the same truncated second order equation \nas in \\cite{Weinberg}.\nIn a medium with nonvanishing pressure one cannot impose the two gauge \nconditions simultaneously (even in the large-scale limit). \nIn Sec. 11.5 of \\cite{Coles} the authors proposed a simple modification\nof the Newtonian theory which again happens to end up with the same incorrect \nequation as in the other books.\nIn the Appendix we elaborate our point.\nFor comments on other related errors in the literature, see Sec. 3.10 in \n\\cite{Ideal} and \\cite{GRG}.\n\nWe would like to remark that these errors found in the synchronous\ngauge are not due to any esoteric aspect of the gauge choice in the \nrelativistic perturbation analyses.\nAlthough not fundamental, the number of errors found in the literature\nconcerning the synchronous gauge seem to indicate the importance of using \nthe proper gauge condition in handling the problems. \n{}For our own choice of the preferred gauge-invariant variables \nsuitable for handling the hydrodynamic perturbation and consequent analyses, \nsee Sec. \\ref{sec:Hydrodynamics}; the reasons for such choices are made in \nSec. \\ref{sec:Six-gauges}.\nWe wish to recall, however, that although one may need to do more algebra \nin tracing the remnant gauge mode, even the synchronous gauge condition\nis adequate for handling many problems as was carefully done in the \noriginal study by Lifshitz in 1946 \\cite{Lifshitz}.\n\nIn this paper we have tried to identify the variables in the relativistic\nperturbation analysis which reproduce the correct Newtonian behavior\nin the pressureless limit.\nIn the first part we have shown that\n$\\delta$ in the comoving gauge ($\\delta_v$) and $v$ and $\\varphi$ \nin the zero-shear gauge ($v_\\chi$ and $\\varphi_\\chi$) \nshow the same behavior as the corresponding Newtonian \nvariables {\\it in general scales}.\nIn fact, these results have already been presented in \\cite{MDE}. \nAlso, various general and asymptotic solutions for every variable\nin the pressureless medium are presented in the Tables of \\cite{MDE}.\nCompared with \\cite{MDE}, in Sec. \\ref{sec:Six-gauges} we tried to reinforce \nthe correspondence by explicitly showing the second order differential \nequations in several gauge conditions.\nWe have also added some additional insights gained after publishing \\cite{MDE}.\nThe second part contains some original results.\nUsing the gauge-invariant variables $\\delta_v$, $v_\\chi$ and $\\varphi_\\chi$ \nwe write the relativistic hydrodynamic cosmological perturbation equations in \nEqs. (\\ref{P-1}-\\ref{P-4}); actually, these equations \nare also known in \\cite{Bardeen-1980}.\nThe new results in this paper are the compact way of deriving the\ngeneral large-scale solutions in Eqs. (\\ref{varphi_chi-sol}-\\ref{v_chi-sol}), \nand the clarification of the general large-scale conservation property of\n$\\varphi_v$ in Eq. (\\ref{varphi_v-sol}).\n\nThe underlying mathematical or physical reasons for \nthe variables $\\delta_v$, $\\varphi_\\chi$, $v_\\chi$, and $\\varphi_v$ \nhaving the distinguished behaviors, compared with many other available\ngauge-invariant combinations, may still deserve further investigation.\n\n\n\\subsection*{Acknowledgments}\n\nWe thank Profs. P. J. E. Peebles and A. M\\'esz\\'aros for useful discussions\nand comments, and Prof. R. Brandenberger for interesting correspondences.\nJH was supported by the KOSEF, Grant No. 95-0702-04-01-3 and\nthrough the SRC program of SNU-CTP.\nHN was supported by the DFG (Germany) and the KOSEF (Korea).\n\n\n\\section*{Appendix: Common errors in the synchronous gauge}\n\nIn the following we correct a minor confusion in the literature\nconcerning perturbation analyses in the synchronous gauge.\nThe argument is based on \\cite{GRG}.\n{}For simplicity, we consider a situation with $K = 0 = \\Lambda$,\n$e = 0 = \\sigma$ and ${\\rm w} = {\\rm constant}$ (thus $c_s^2 = {\\rm w}$).\n\nThe equation for the density perturbation in the {\\it comoving gauge}\nis given in Eq. (\\ref{delta-eq}).\nIn our case we have\n$$\n \\ddot \\delta_v + ( 2- 3 {\\rm w}) H \\dot \\delta_v\n + \\left[ c_s^2 {k^2 \\over a^2}\n - 4 \\pi G \\mu (1 - {\\rm w}) ( 1 + 3 {\\rm w} ) \\right] \\delta_v = 0.\n \\eqno{(A1)}\n$$\nThe solution in the super-sound-horizon scale is presented in\nEq. (\\ref{Sol-w=const}).\n\nIn the {\\it synchronous gauge}, from Eqs. (\\ref{eq5}-\\ref{eq7}) we have\n(the subindex $\\alpha$ indicates the synchronous gauge variable)\n$$\n \\ddot \\delta_\\alpha + 2 H \\dot \\delta_\\alpha\n + \\left[ c_s^2 {k^2 \\over a^2}\n - 4 \\pi G \\mu ( 1+ {\\rm w} ) ( 1 + 3 {\\rm w} ) \\right] \\delta_\\alpha\n + 3 {\\rm w} ( 1 + {\\rm w}) {k \\over a} H v_\\alpha = 0,\n \\eqno{(A2)}\n$$\n$$\n \\dot v_\\alpha + ( 1- 3 {\\rm w}) H v_\\alpha\n - {k \\over a} {{\\rm w} \\over 1 + {\\rm w}} \\delta_\\alpha = 0.\n \\eqno{(A3)}\n$$\nNotice that for ${\\rm w} \\neq 0$ these two equations are generally coupled\neven in the large-scale limit.\n{}From these two equations we can derive a third order differential equation\nfor $\\delta_\\alpha$ as\n$$\n \\delta^{\\cdot\\cdot\\cdot}_\\alpha\n + {11 - 3 {\\rm w} \\over 2} H \\ddot \\delta_\\alpha\n + \\left( {5 - 24 {\\rm w} - 9 {\\rm w}^2 \\over 2} H^2\n + c_s^2 {k^2 \\over a^2} \\right) \\dot \\delta_\\alpha\n$$\n$$\n + \\left[ {3 \\over 4} \\left( 1 + {\\rm w} \\right)\n \\left( 1 + 3 {\\rm w} \\right) \\left( -1 + 9 {\\rm w} \\right) H^3\n + {3 \\over 2} \\left( 1 + {\\rm w} \\right) H c_s^2 {k^2 \\over a^2}\n \\right] \\delta_\\alpha = 0.\n \\eqno{(A4)}\n$$\nIn the large-scale limit the solutions are\n$$\n \\delta_\\alpha \\quad \\propto \\quad\n t^{{ 2(1+ 3{\\rm w}) \\over 3(1 + {\\rm w})}}, \\quad\n t^{{9{\\rm w} - 1 \\over 3( 1+ {\\rm w})}}, \\quad\n t^{-1}.\n \\eqno{(A5)}\n$$\nIn Appendix B of \\cite{GRG} we showed that\nthe third solution with $\\delta_\\alpha \\propto t^{-1}$\nis nothing but a gauge mode for a medium with ${\\rm w} \\neq 0$.\n[As mentioned before, the combination $\\delta_\\alpha$ is not gauge-invariant.\n {}From Eq. (\\ref{GT}) we have\n $\\delta_\\alpha \\equiv \\delta + 3 H ( 1 + {\\rm w}) \\int^t \\alpha dt$\n and the lower bound of the integration gives rise to the gauge mode \n which is proportional to $t^{-1}$.]\n{}For a pressureless case the physical decaying solution also behaves \nas $t^{-1}$ and the second solution in Eq. (A5) is invalid,\nsee Sec. 4 in \\cite{GRG}.\nIn Eq. (B5) of \\cite{GRG} we derived the relation between solutions\nin the two gauges explicitly; the growing solutions are the same in both gauges\nwhereas the decaying solutions differ by a factor\n$(k \/ aH)^2 \\propto t^{2(1 + 3 {\\rm w}) \\over 3 (1 + {\\rm w})}$.\n\nNow, we would like to point out that, in a medium with ${\\rm w} \\neq 0$,\none cannot ignore the last term in Eq. (A2)\neven in the large-scale limit.\nIf we {\\it incorrectly} ignore the last term in Eq. (A2) we recover\nthe {\\it wrong equation} in the textbooks mentioned in \nSec. \\ref{sec:Discussion} which is\n$$\n \\ddot \\delta_\\alpha + 2 H \\dot \\delta_\\alpha\n + \\left[ c_s^2 {k^2 \\over a^2}\n - 4 \\pi G \\mu ( 1 + {\\rm w} ) ( 1 + 3 {\\rm w} ) \\right] \\delta_\\alpha\n = 0.\n \\eqno{(A6)}\n$$\nIn the large-scale we have solutions\n$$\n \\delta_\\alpha \\quad \\propto \\quad\n t^{{ 2(1+ 3{\\rm w}) \\over 3(1 + {\\rm w})}}, \\quad\n t^{-1}.\n \\eqno{(A7)}\n$$\nBy ignoring the last term in Eq. (A2), in the large-scale\nlimit we happen to recover the fictitious gauge mode under the price of losing\nthe physical decaying solution [the second solution in Eq. (A5)].\nThus, in the large-scale limit one cannot ignore the last term in Eq. (A2)\nin a medium with a general pressure; the reason is obvious if we see\nEq. (A4).\nAlso, one cannot impose both the synchronous gauge condition and the\ncomoving gauge condition simultaneously.\nIf we simultaneously impose such two conditions,\nthus setting $v \\equiv 0 \\equiv \\alpha$, from Eq. (\\ref{eq7}) we have\n$$\n 0 = {k \\over a \\left( 1 + {\\rm w} \\right) } \\left( c_s^2 \\delta\n + {e \\over \\mu} - {2\\over 3} {k^2 - 3K \\over a^2} {\\sigma \\over \\mu}\n \\right).\n \\eqno{(A8)}\n$$\nThus, for $e = 0 = \\sigma$ and medium with the nonvanishing pressure\nwe have $\\delta = 0$ which is a meaningless system;\nthis argument remains valid even in the large-scale limit.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{ Introduction }\nIt has long been a problem that the $\\Delta I=3\/2$ \n$K^+\\to\\pi^+\\pi^0$ amplitude calculated in quenched lattice QCD is \nabout a factor two too large compared to the experimental value\\cite{BSReview}.\nIn this article we report a new study of this problem,\nincorporating various theoretical and technical advances \nin recent years for analysis.\nIn particular we discuss in detail how one-loop corrections of \nchiral perturbation theory (CHPT) recently calculated\\cite{ONECPTH} affect\nphysical predictions for the decay amplitude from lattice QCD simulations. \n\nOur simulation is carried out in quenched lattice QCD employing \nthe standard plaquette action for gluons at $\\beta=6.1$ and the Wilson action for quarks. \nTwo lattice sizes, $24^3\\times 64$ and $32^3\\times 64$, are employed.\nWe take up, down and strange quarks to be degenerate,\nand make measurements at four values of the common hopping parameter,\n$\\kappa = 0.1520$, $0.1530$, $0.1540$ and $0.1543$, which correspond to\n$m_\\pi\/m_\\rho = 0.797$, $0.734$, $0.586$ and $0.515$.\n\\section{ Extraction of decay amplitude }\nThe 4-quark operator most relevant for the \n$\\Delta I=3\/2$ $K\\to\\pi\\pi$ decay is\n$Q_+ = [ ( \\bar{s} d )_L\n ( \\bar{u} u )_L\n + ( \\bar{s} u )_L\n ( \\bar{u} d )_L ] \/ 2$.\nWe extract the decay amplitude from the 4-point correlation function defined by\n$M_Q = \\langle 0 | W_+ W_0 Q_{+}(t) W_K | 0 \\rangle$ \nwhere $W_{0,+,K}$ are wall sources for $\\pi^0$, $\\pi^+$ and $K^+$. \nIn our calculations on a lattice of a temporal size $T=64$, \nthe walls are placed at the time slices $t_{K^+}=4$, $t_{\\pi^+}=59$ and $t_{\\pi^0}=60$.\nThe mesons are all created at rest, and the 4-quark operator $Q_+$ is \nprojected to zero spatial momentum.\n\nWe can use two types of factors to remove the normalization factors in $M_Q$.\nIf we define\n\\begin{eqnarray}\n\\lefteqn{ M_W =\n \\langle 0 | W_0 \\pi^{0} ( t ) | 0 \\rangle\n \\langle 0 | W_+ \\pi^{+} ( t ) | 0 \\rangle\n \\langle 0 | K(t) W_K | 0 \\rangle ,} && \\nonumber \\\\\n\\lefteqn{ M_P = \n \\langle 0 | K(t) W_K | 0 \\rangle \n \\langle 0 | W_+ W_0 \\pi^{+} ( t ) \\pi^{0} ( t ) | 0 \\rangle , } && \n\\end{eqnarray}\nwe find\n$R_W \\equiv M_Q \/ M_W$\n$\\sim \\langle \\pi^+\\pi^0 | Q_+ | K^+ \\rangle$\n$ \/ \\langle \\pi | \\pi | 0 \\rangle^3 \\cdot {\\rm exp}(t-t_+)\\Delta$\nand\n$R_P \\equiv M_Q \/ M_P$ \n$\\sim \\langle \\pi^+\\pi^0 | Q_+ | K^+ \\rangle \/ \\langle \\pi | \\pi | 0 \\rangle^3$\nfor $t_K\\ll t\\ll t_+, t_0$, where $\\Delta=m_{\\pi\\pi}-2m_\\pi$ is the \nmass shift of the 2-pion state due to finite lattice size \neffects~\\cite{LUSHER}.\n \nIn Fig.~\\ref{RWFRPF} we plot $\\langle \\pi | \\pi | 0 \\rangle^3 \\cdot R_W$ and \n$\\langle \\pi | \\pi | 0\\rangle^3 \\cdot R_P$ at \n$\\kappa=0.1530$ as a function of time $t$ of the weak operator.\nWe clearly observe a non-vanishing slope for $R_W$, while $R_P$ exhibits a \nplateau as expected.\nThe decay amplitude can be obtained by fitting $R_W$ to a single exponential \nor $R_P$ to a constant. We find the results to be mutually consistent \nwithin the statistical error.\n\\begin{figure}[t]\n\\vspace*{-0.2cm}\n\\centerline{\\epsfxsize=7.0cm \\epsfbox{L24L32all.1.eps}}\n\\vspace*{-1cm}\n\\caption{\\label{RWFRPF}\n$R_W \\cdot \\langle \\pi | \\pi | 0 \\rangle^3$ and $R_P\\langle \\pi | \\pi | 0 \\rangle^3$\nat $\\kappa=0.153$.\nOpen and filled circles refer to data for $24^3$ and $32^3$ lattices.\n}\n\\vspace*{-0.7cm}\n\\end{figure}\n\\begin{figure}[t]\n\\vspace{-0.6cm}\n\\centerline{\\epsfxsize=7.0cm \\epsfbox{o00210o11210.00.eps}}\n\\vspace{-1.0cm}\n\\caption{ \\label{COMPARISON}\nComparison of our results normalized by the experimental value obtained \nwith tree-level CHPT relation for $q^*=\\pi\/a$ at $\\beta=6.1$ \nwith those of previous work{\\protect\\cite{BSKpipi}} at $\\beta=6.0$ (crosses).\nOpen symbols refer to results with the traditional normalization \n$\\protect\\sqrt{2\\kappa}$ and no tadpole improvement, \nand filled ones are those with the KLM normalization and tadpole improvement.\n}\n\\vspace{-0.7cm}\n\\end{figure}\n\nThe one-loop renormalization factor relating the operator $Q_+$ on the lattice \nand that in the continuum was obtained in Refs.~\\cite{Z-M,Z-BS}. \nWe set $q^*=1\/a$ or $\\pi\/a$ as the matching point and employ tadpole-improved\nperturbation theory with \nthe mean-field improved $\\overline{\\rm MS}$ coupling constant\\cite{TPI}. \nQuark fields are normalized with the KLM factor~\\cite{KLM}, \n$\\sqrt{1-3\\kappa\/4\\kappa_c}$.\n\\section{ Result }\nIn earlier calculations tree-level chiral perturbation theory (CHPT) \nformula was used to obtain the physical amplitude $A^{phys}$ from the \nlattice counterpart $A^{lat}$. The formula takes the form \n$ A^{phys} = ( m_K^2 - m_\\pi^2 )\/( 2 M_\\pi^2 ) \\cdot A^{lat}$\nwhere $m_K=497 {\\rm MeV}$ and $m_\\pi=136 {\\rm MeV}$ are physical masses, \nand $M_\\pi$ is the degenerate mass of $K$ and $\\pi$ mesons on the lattice.\nClearly the physical amplitude should be independent of the \nlattice mass $M_\\pi$ if the matching formula is exact.\n\nIn Fig.~\\ref{COMPARISON} we compare results for the physical decay amplitude \nfrom a previous work\\cite{BSKpipi} carried out at $\\beta=6.0$ on a $24^3$ \nspatial lattice with those of our simulation at $\\beta=6.1$. \nOur values obtained with the conventional quark normalization $\\sqrt{2\\kappa}$\n(crosses and open squares) are consistent with their results,\nbut are larger than the experimental result roughly by a factor two.\n\nLet us also note with our results (circles and squares) \nthat (i) the KLM normalization and tadpole improvement of \nthe operator have a significant effect on the amplitude, (ii) there is \na clear dependence of the amplitude on the lattice meson mass $M_\\pi$, \nand (iii) a significant finite-size effect is observed between the two \nlattice sizes.\nThese features show that the tree-level CHPT is inadequate to extract \nthe physical amplitude from lattice results.\n\nIn Fig.~\\ref{FRHO} we show how predictions for the physical amplitude \nchange if we apply the one-loop formula of CHPT~\\cite{ONECPTH} to our results.\nHere $\\Lambda^q$ and $\\Lambda^{cont}$ are the cutoffs of CHPT for quenched and full theory.\nIn converting to physical values we ignore effects of $O(p^4)$ contact terms \nof the CHPT Lagrangian since their values are not well known.\nAn interesting point is that a size dependence seen with \nthe tree-level analysis in Fig.~\\ref{COMPARISON} is absent.\nAnother important point is that the magnitude of the amplitude decreases \nby $30-40$\\% over the range of meson mass covered in our simulation.\nThe amplitude depends significantly on the quenched lattice cutoff \n$\\Lambda^q$, in particular toward larger values of $M_\\pi$.\n\n\\begin{figure}\n\\vspace*{-0.6cm}\n\\centerline{\\epsfxsize=7.0cm \\epsfbox{o11210.1113.eps}}\n\\vspace*{-1cm}\n\\caption{ \\label{FRHO} \nDecay amplitude normalized by experimental value for $q^*=\\pi\/a$ \nobtained with one-loop CHPT for $\\Lambda^{cont} = 770{\\rm MeV}$.\nCircles and squares refer to data for $24^3$ and $32^3$ spatial sizes. \nOpen symbols are for $\\Lambda^{q}=770{\\rm MeV}$ and filled ones for \n$1\\ {\\rm GeV}$.\n}\n\\vspace*{-7mm}\n\\end{figure}\nOur results in Fig.~\\ref{FRHO} show that a sizable meson mass dependence still remains.\nThis may be attributed to $O(p^4)$ contact terms which were not taken into account.\nIf we assume that the contact terms of the full theory are as small \nas suggested by the phenomenological analysis~\\cite{BK},\nwe can estimate the physical amplitude by a linear extrapolation of $M_\\pi^2$ to the chiral limit.\nIn Table~\\ref{FINALFIT} we list results for several choices of the cutoff\nand the operator matching point $q^*$.\nAs is expected from Fig.~\\ref{FRHO}, variation of results on the \nchoice of the cutoff parameters is small, being of the order of 10\\%.\nAllowing for this uncertainty, the values in Table~\\ref{FINALFIT}\nare consistent with the experimental value of $10.4\\times 10^{-3} {\\rm GeV}^3$.\n\\section{ Conclusions }\nThe present study has shown that the quenched result for the \n$K^+\\to\\pi^+ \\pi^0 $ decay amplitude agrees with experiment much better \nthan previously thought, especially if the chiral limit is taken \nfor the lattice meson mass.\nAway from the chiral limit, the amplitude depends significantly on the cutoff scales of CHPT, however. \nThis problem is largely ascribed to a mismatch of chiral logarithms \nbetween the quenched and full theories\\cite{ONECPTH}. \nIn order to obtain $K^+\\to\\pi^+ \\pi^0 $ decay amplitude free from this uncertainty,\nwe need simulations in full QCD.\n\n\\begin{table}[t]\n\\setlength{\\tabcolsep}{0.2pc}\n\\begin{tabular}{llllll}\n\\hline\n\\hline\n&&\\multicolumn{4}{c}{$C_+\\langle \\pi\\pi |Q_+| K \\rangle (\\times 10^{-3} {\\rm GeV}^3)$}\\\\\n$\\Lambda^{cont}$ & $\\Lambda^{q}$ & \\multicolumn{2}{c}{$24^3$}&\\multicolumn{2}{c}{$32^3 $}\\\\\n({\\small GeV}) & ({\\small GeV}) &$q^*=\\frac{1}{a}$&$q^*=\\frac{\\pi}{a}$&$q^*=\\frac{1}{a}$&$q^*=\\frac{\\pi}{a}$\\\\[1mm]\n\\hline\n0.77 & 0.77 & $ 9.3(19)$&$10.2(21)$ & $8.9(17)$&$ 9.7(19)$ \\\\\n0.77 & 1.0 & $ 9.4(13)$&$10.3(14)$ & $8.8(11)$&$9.6(12)$ \\\\\n1.0 & 0.77 & $10.3(21)$&$11.3(23)$ & $9.8(19)$&$10.7(21)$ \\\\\n1.0 & 1.0 & $10.4(14)$&$11.4(15)$ & $9.7(12)$&$10.6(13)$ \\\\\n\\hline\n\\hline\n\\end{tabular}\n\\vspace*{3mm}\n\\caption{\\label{FINALFIT}\nResults of linear extrapolation of the decay amplitude to $M_\\pi^2=0$. \nStatistical and extrapolation errors are combined.\nThe experimental value is $10.4\\times 10^{-3} {\\rm GeV}^3$.\n}\n\\vspace*{-7mm}\n\\end{table}\n\\hfill\\break\nThis work is supported by the Supercomputer \nProject (No.~97-15) of High Energy Accelerator Research Organization (KEK),\nand also in part by the Grants-in-Aid of \nthe Ministry of Education (Nos. 08640349, 08640350, 08640404,\n09246206, 09304029, 09740226).\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzmfmx b/data_all_eng_slimpj/shuffled/split2/finalzzmfmx new file mode 100644 index 0000000000000000000000000000000000000000..7022f5dd0e3081924d4097e3a60f495d53780819 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzmfmx @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\n\nIn 1932, the physicist P. Jordan proposed a program to discover a new algebraic setting for quantum mechanics, so Jordan algebras were created in this proceeding. Moreover, Jordan algebras were truned out to have illuminating connections with many areas of mathematics. A. A. Albert renamed them Jordan algebras and developed a successful structure theory for Jordan algebras over any field of characteristic zero in \\cite{A1}.\n\nHom-type algebras have been studied by many authors in \\cite{A2}, \\cite{A3}, \\cite{L1}, \\cite{D1}. In general, Hom-type algebras are a kind of algebras which are obtained by twisting the identity defining the structure with one or several linear maps(called the twisted maps) on the basis of the original algebras. The concept of Hom-Jordan algebras was first introduced by A. Makhlouf in his paper \\cite{M2}. He introduced a Hom-Jordan algebra and showed that it fits with the Hom-associative structure, that is a Hom-associative algebra leads to a Hom-Jordan algebra by considering a plus algebra.\n\nAs is well known, derivations and generalized derivations play an important role in the research of structure and property in algebraic system. The research on generalized derivations was started by Leger and Luks in \\cite{L2}. They gave many properties about generalized derivations of Lie algebras. From then on, many authors generalized the results in \\cite{L2} to other algebras. R. X. Zhang and Y. Z. Zhang generalized it to Lie superalgebras in \\cite{Z1}. And L. Y. Chen, Y. Ma and L. Ni generalized the results to Lie color algebras successfully in \\cite{C1}. J. Zhou, Y. J. Niu and L. Y. chen generalized the above results to Hom-Lie algebras in \\cite{Z2}. As for Jordan algebras, A. I. Shestakov generalized the results into Jordan superalgebras in \\cite{S1}. He showed that for semi-simple Jordan algebras over any field of characteristic $0$ or simple Jordan algebras over any field of characteristic other than $2$, the generalized derivations are the sum of the derivations and the centroids with some exceptions.\n\nThe purpose of this paper is to generalized some beautiful results in \\cite{Z2} and \\cite{M3} to Hom-Jordan algebras. We proceed as follow. In Section \\ref{se:2}, we recall some basic definitions and propositions which will be used in what follows. In Section \\ref{se:3}, we'll give some properties about generalized derivations of Hom-Jordan algebras and their Hom-subalgebras. In Section \\ref{se:4}, we show that quasiderivations of a Hom-Jordan algebra can be embedded as derivations into a larger Hom-Jordan algebra and obtain a direct sum decomposition of $Der(\\breve{V})$ when the centralizer of $V$ equals to zero. In Section \\ref{se:5}, we give some propositions with respect to the centroids of Hom-Jordan algebras.\n\n\\section{Preliminaries}\\label{se:2}\n\n\\begin{defn}\\cite{M1}\nAn algebra $J$ over a field $\\rm{F}$ is a Jordan algebra satisfying for any $x, y \\in J$,\n\\begin{enumerate}[(1)]\n\\item $x \\circ y = y \\circ x$;\n\\item $(x^{2} \\circ y) \\circ x = x^{2} \\circ (y \\circ x)$.\n\\end{enumerate}\n\\end{defn}\n\\begin{defn}\\cite{M2}\nA Hom-Jordan algebra over a field $\\rm{F}$ is a triple $(V, \\mu, \\alpha)$ consisting of a linear space $V$, a bilinear map $\\mu : V \\times V \\rightarrow V$ which is commutative and a linear map $\\alpha : V \\rightarrow V$ satisfying for any $x, y \\in V$,\n\\[\\mu(\\alpha^{2}(x), \\mu(y, \\mu(x, x))) = \\mu(\\mu(\\alpha(x), y), \\alpha(\\mu(x, x))),\\]\nwhere $\\alpha^{2} = \\alpha \\circ \\alpha$.\n\\end{defn}\n\\begin{defn}\nA Hom-Jordan algebra $(V, \\mu, \\alpha)$ is called multiplicative if for any $x, y \\in V$, $\\alpha(\\mu(x, y)) = \\mu(\\alpha(x), \\alpha(y))$.\n\\end{defn}\n\\begin{defn}\\cite{B1}\nA subspace $W \\subseteq V$ is a Hom-subalgebra of $(V, \\mu, \\alpha)$ if $\\alpha(W) \\subseteq W$ and\n\\[\\mu(x, y) \\in W,\\quad\\forall x, y \\in W.\\]\n\\end{defn}\n\\begin{defn}\\cite{B1}\nA subspace $W \\subseteq V$ is a Hom-ideal of $(V, \\mu, \\alpha)$ if $\\alpha(W) \\subseteq W$ and\n\\[\\mu(x, y) \\in W,\\quad\\forall x \\in W, y \\in V.\\]\n\\end{defn}\n\\begin{defn}\\cite{B1}\nLet $(V, \\mu, \\alpha)$ and $(V^{'}, \\mu^{'}, \\beta)$ be two Hom-Jordan algebras. A linear map $\\phi : V \\rightarrow V^{'}$ is said to be a homomorphism of Hom-Jordan algebras if\n\\begin{enumerate}[(1)]\n\\item $\\phi(\\mu(x, y)) = \\mu^{'}(\\phi(x), \\phi(y))$;\n\\item $\\phi \\circ \\alpha = \\beta \\circ \\phi$.\n\\end{enumerate}\n\\end{defn}\n\\begin{lem}\\cite{B1}\\label{le:2.7}\nLet $(V, \\mu, \\alpha)$ be a Hom-Jordan algebra over a field $\\rm{F}$, we define a subspace $\\mathcal{W}$ of $End(V)$ where $\\mathcal{W} = \\{w \\in End(V) | w \\circ \\alpha = \\alpha \\circ w\\}$, $\\sigma : \\mathcal{W} \\rightarrow \\mathcal{W}$ is a linear map satisfying $\\sigma(w) = \\alpha \\circ w$.\n\\begin{enumerate}[(1)]\n\\item A triple $(\\mathcal{W}, \\nu, \\sigma)$, where the multiplication $\\nu : \\mathcal{W} \\times \\mathcal{W} \\rightarrow \\mathcal{W}$ is defined for $w_{1}, w_{2} \\in \\mathcal{W}$ by\n \\[\\nu(w_{1}, w_{2}) = w_{1} \\circ w_{2} + w_{2} \\circ w_{1},\\]\n is a Hom-Jordan algebra over $\\rm{F}$.\n\\item A triple $(\\mathcal{W}, \\nu^{'}, \\sigma)$, where the multiplication $\\nu^{'} : \\mathcal{W} \\times \\mathcal{W} \\rightarrow \\mathcal{W}$ is defined for $w_{1}, w_{2} \\in \\mathcal{W}$ by\n \\[\\nu^{'}(w_{1}, w_{2}) = w_{1} \\circ w_{2} - w_{2} \\circ w_{1},\\]\n is a Hom-Lie algebra over $\\rm{F}$.\n\\end{enumerate}\n\\end{lem}\n\\begin{defn}\\cite{B1}\nFor any nonnegative integer $k$, a linear map $D : V \\rightarrow V$ is called an $\\alpha^{k}$-derivation of the Hom-Jordan algebra $(V, \\mu, \\alpha)$, if\n\\begin{enumerate}[(1)]\n\\item $D \\circ \\alpha = \\alpha \\circ D$;\n\\item $D(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D(y)),\\quad\\forall x, y \\in V$.\n\\end{enumerate}\n\\end{defn}\n\\begin{lem}\\cite{B1}\\label{le:2.9}\nFor any $D \\in Der_{\\alpha^{k}}(V)$ and $D^{'} \\in Der_{\\alpha^{s}}(V)$, define their commutator $\\nu^{'}(D, D^{'})$ as usual:\n\\[\\nu^{'}(D, D^{'}) = D \\circ D^{'} - D^{'} \\circ D.\\]\nThen $\\nu^{'}(D, D^{'}) \\in Der_{\\alpha^{k + s}}(V)$.\n\\end{lem}\n\nLet $(V, \\mu, \\alpha)$ be a multiplicative Hom-Jordan algebra. We denote the set of all $\\alpha^{k}$-derivations by $Der_{\\alpha^{k}}(V)$, then $Der(V) = \\dotplus_{k \\geq 0}Der_{\\alpha^{k}}(V)$ provided with the commutator and the following map\n\\[\\sigma : Der(V) \\rightarrow Der(V);\\quad \\sigma(D) = \\alpha \\circ D\\]\nis a Hom-subalgebra of $(\\mathcal{W}, \\nu^{'}, \\sigma)$ according to Lemma \\ref{le:2.9} and is called the derivation algebra of $(V, \\mu, \\alpha)$.\n\\begin{defn}\nFor any nonnegative integer $k$, a linear map $D \\in End(V)$ is said to be a generalized $\\alpha^{k}$-derivation of $(V, \\mu, \\alpha)$, if there exist two linear maps $D^{'}, D^{''} \\in End(V)$ such that\n\\begin{enumerate}[(1)]\n\\item $D \\circ \\alpha = \\alpha \\circ D,\\; D^{'} \\circ \\alpha = \\alpha \\circ D^{'},\\; D^{''} \\circ \\alpha = \\alpha \\circ D^{''}$;\n\\item $\\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D^{'}(y)) = D^{''}(\\mu(x, y)),\\quad\\forall x, y \\in V$.\n\\end{enumerate}\n\\end{defn}\n\\begin{defn}\\label{defn:2.11}\nFor any nonnegative integer $k$, a linear map $D \\in End(V)$ is said to be an $\\alpha^{k}$-quasiderivation of $(V, \\mu, \\alpha)$, if there exist a linear map $D^{'}\\in End(V)$ such that\n\\begin{enumerate}[(1)]\n\\item $D \\circ \\alpha = \\alpha \\circ D,\\; D^{'} \\circ \\alpha = \\alpha \\circ D^{'}$;\n\\item $\\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D(y)) = D^{'}(\\mu(x, y)),\\quad\\forall x, y \\in V$.\n\\end{enumerate}\n\\end{defn}\n\nLet $GDer_{\\alpha^{k}}(V)$ and $QDer_{\\alpha^{k}}(V)$ be the sets of generalized $\\alpha^{k}$-derivations and of $\\alpha^{k}$-quasiderivations, respectively. That is\n\\[GDer(V) = \\dotplus_{k \\geq 0}GDer_{\\alpha^{k}}(V),\\quad QDer(V) = \\dotplus_{k \\geq 0}QDer_{\\alpha^{k}}(V).\\]\nIt's easy to verify that both $GDer(V)$ and $QDer(V)$ are Hom-subalgebras of $(\\mathcal{W}, \\nu^{'}, \\sigma)$(See Proposition \\ref{prop:3.1}).\n\\begin{defn}\nFor any nonnegative integer $k$, a linear map $D \\in End(V)$ is said to be an $\\alpha^{k}$-centroid of $(V, \\mu, \\alpha)$, if\n\\begin{enumerate}[(1)]\n\\item $D \\circ \\alpha = \\alpha \\circ D$;\n\\item $\\mu(D(x), \\alpha^{k}(y)) = \\mu(\\alpha^{k}(x), D(y)) = D(\\mu(x, y)),\\quad\\forall x, y \\in V$.\n\\end{enumerate}\n\\end{defn}\n\\begin{defn}\nFor any nonnegative integer $k$, a linear map $D \\in End(V)$ is said to be an $\\alpha^{k}$-quasicentroid of $(V, \\mu, \\alpha)$, if\n\\begin{enumerate}[(1)]\n\\item $D \\circ \\alpha = \\alpha \\circ D$;\n\\item $\\mu(D(x), \\alpha^{k}(y)) = \\mu(\\alpha^{k}(x), D(y)),\\quad\\forall x, y \\in V$.\n\\end{enumerate}\n\\end{defn}\nLet $C_{\\alpha^{k}}(V)$ and $QC_{\\alpha^{k}}(V)$ be the sets of $\\alpha^{k}$-centroids and of $\\alpha^{k}$-quasicentroids, respectively. That is\n\\[C(V) = \\dotplus_{k \\geq 0}C_{\\alpha^{k}}(V),\\quad QC(V) = \\dotplus_{k \\geq 0}QC_{\\alpha^{k}}(V).\\]\n\\begin{defn}\nFor any nonnegative integer $k$, a linear map $D \\in End(V)$ is said to be an $\\alpha^{k}$-central derivation of $(V, \\mu, \\alpha)$, if\n\\begin{enumerate}[(1)]\n\\item $D \\circ \\alpha = \\alpha \\circ D$;\n\\item $\\mu(D(x), \\alpha^{k}(y)) = D(\\mu(x, y)) = 0,\\quad\\forall x, y \\in V$.\n\\end{enumerate}\n\\end{defn}\n\nWe denote the set of all $\\alpha^{k}$-central derivations by $ZDer_{\\alpha^{k}}(V)$, then $ZDer(V) = \\dotplus_{k \\geq 0}ZDer_{\\alpha^{k}}(V)$.\n\nAccording to the definitions, it's easy to show that $ZDer(V) \\subseteq Der(V) \\subseteq QDer(V) \\subseteq GDer(V) \\subseteq End(V)$.\n\\begin{defn}\nLet $(V, \\mu, \\alpha)$ be a Hom-Jordan algebra. If $Z(V) = \\{x \\in V | \\mu(x, y) = 0,\\quad\\forall y \\in V\\}$, then $Z(V)$ is called the centralizer of $(V, \\mu, \\alpha)$.\n\\end{defn}\n\n\\section{Generalized derivation algebras and their subalgebras}\\label{se:3}\n\nAt first, we give some basic properties of center derivation algebras, quasiderivation algebras and generalized derivation algebras of a Hom-Jordan algebra.\n\\begin{prop}\\label{prop:3.1}\nSuppose that $(V, \\mu, \\alpha)$ is a multiplicative Hom-Jordan algebra. Then the following statements hold:\n\\begin{enumerate}[(1)]\n\\item $GDer(V)$, $QDer(V)$ and $C(V)$ are Hom-subalgebras of $(\\mathcal{W}, \\nu^{'}, \\sigma)$;\n\\item $ZDer(V)$ is a Hom-ideal of $Der(V)$.\n\\end{enumerate}\n\\end{prop}\n\\begin{proof}\n(1). Suppose that $D_{1} \\in GDer_{\\alpha^{k}}(V)$, $D_{2} \\in GDer_{\\alpha^{s}}(V)$. Then for any $x, y \\in V$,\n\\begin{align*}\n&\\mu(\\sigma(D_{1})(x), \\alpha^{k + 1}(y)) = \\mu(\\alpha \\circ D_{1}(x), \\alpha^{k + 1}(y)) = \\alpha(\\mu(D_{1}(x), \\alpha^{k}(y)))\\\\\n&= \\alpha(D_{1}^{''}(\\mu(x, y)) - \\mu(\\alpha^{k}(x), D_{1}^{'}(y))) = \\sigma(D_{1}^{''})(\\mu(x, y)) - \\mu(\\alpha^{k + 1}(x), \\sigma(D_{1}^{'})(y)).\n\\end{align*}\nSince $\\sigma(D_{1}^{''}), \\sigma(D_{1}^{'}) \\in End(V)$, we have $\\sigma(D_{1}) \\in GDer_{\\alpha^{k + 1}}(V)$.\n\\begin{align*}\n&\\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y))\\\\\n&= \\mu(D_{1} \\circ D_{2}(x), \\alpha^{k + s}(y)) - \\mu(D_{2} \\circ D_{1}(x), \\alpha^{k + s}(y))\\\\\n&= D_{1}^{''}(\\mu(D_{2}(x), \\alpha^{s}(y))) - \\mu(\\alpha^{k}(D_{2}(x)), D_{1}^{'}(\\alpha^{s}(y))) - D_{2}^{''}(\\mu(D_{1}(x), \\alpha^{k}(y))) +\\\\ &\\mu(\\alpha^{s}(D_{1}(x)), D_{2}^{'}(\\alpha^{k}(y)))\\\\\n&= D_{1}^{''}(D_{2}^{''}(\\mu(x, y)) - \\mu(\\alpha^{s}(x), D_{2}^{'}(y))) - \\mu(\\alpha^{k}(D_{2}(x)), D_{1}^{'}(\\alpha^{s}(y)))\\\\\n&- D_{2}^{''}(D_{1}^{''}(\\mu(x, y)) - \\mu(\\alpha^{k}(x), D_{1}^{'}(y))) + \\mu(\\alpha^{s}(D_{1}(x)), D_{2}^{'}(\\alpha^{k}(y)))\\\\\n&= D_{1}^{''} \\circ D_{2}^{''}(\\mu(x, y)) - D_{1}^{''}(\\mu(\\alpha^{s}(x), D_{2}^{'}(y))) - \\mu(\\alpha^{k}(D_{2}(x)), D_{1}^{'}(\\alpha^{s}(y)))\\\\\n&- D_{2}^{''} \\circ D_{1}^{''}(\\mu(x, y)) + D_{2}^{''}(\\mu(\\alpha^{k}(x), D_{1}^{'}(y))) + \\mu(\\alpha^{s}(D_{1}(x)), D_{2}^{'}(\\alpha^{k}(y)))\\\\\n&= D_{1}^{''} \\circ D_{2}^{''}(\\mu(x, y)) - \\mu(D_{1}(\\alpha^{s}(x)), \\alpha^{k}(D_{2}^{'}(y))) - \\mu(\\alpha^{k + s}(x), D_{1}^{'}(D_{2}^{'}(y)))\\\\\n&- \\mu(\\alpha^{k}(D_{2}(x)), D_{1}^{'}(\\alpha^{s}(y))) - D_{2}^{''} \\circ D_{1}^{''}(\\mu(x, y)) + \\mu(D_{2}(\\alpha^{k}(x)), \\alpha^{s}(D_{1}^{'}(y)))\\\\\n&+ \\mu(\\alpha^{k + s}(x), D_{2}^{'}(D_{1}^{'}(y))) + \\mu(\\alpha^{s}(D_{1}(x)), D_{2}^{'}(\\alpha^{k}(y)))\\\\\n&= D_{1}^{''} \\circ D_{2}^{''}(\\mu(x, y)) - D_{2}^{''} \\circ D_{1}^{''}(\\mu(x, y)) - \\mu(\\alpha^{k + s}(x), D_{1}^{'}(D_{2}^{'}(y)))\\\\\n&+ \\mu(\\alpha^{k + s}(x), D_{2}^{'}(D_{1}^{'}(y)))\\\\\n&= \\nu^{'}(D_{1}^{''}, D_{2}^{''})(\\mu(x, y)) - \\mu(\\alpha^{k + s}(x), \\nu^{'}(D_{1}^{'}, D_{2}^{'})(y)).\n\\end{align*}\nSince $\\nu^{'}(D_{1}^{''}, D_{2}^{''}), \\nu^{'}(D_{1}^{'}, D_{2}^{'}) \\in End(V)$, we have $\\nu^{'}(D_{1}, D_{2}) \\in GDer_{\\alpha^{k + s}}(V)$.\n\nTherefore, $GDer(V)$ is a Hom-subalgebra of $(\\mathcal{W}, \\nu^{'}, \\sigma)$.\n\nSimilarly, we have $QDer(V)$ is a Hom-subalgebra of $(\\mathcal{W}, \\nu^{'}, \\sigma)$.\n\nSuppose that $D_{1} \\in C_{\\alpha^{k}}(V)$, $D_{2} \\in C_{\\alpha^{s}}(V)$. Then for any $x, y \\in V$,\n\\[\\sigma(D_{1})(\\mu(x, y)) = \\alpha \\circ D_{1}(\\mu(x, y)) = \\alpha(\\mu(D_{1}(x), \\alpha^{k}(y))) = \\mu(\\sigma(D_{1})(x), \\alpha^{k + 1}(y)).\\]\nSimilarly, we have $\\sigma(D_{1})(\\mu(x, y)) = \\mu(\\alpha^{k + 1}(x), \\sigma(D_{1})(y))$. Hence, we have $\\sigma(D_{1}) \\in C_{\\alpha^{k + 1}}(V)$.\n\\begin{align*}\n&\\nu^{'}(D_{1}, D_{2})(\\mu(x, y))\\\\\n&= D_{1} \\circ D_{2}(\\mu(x, y)) - D_{2} \\circ D_{1}(\\mu(x, y))\\\\\n&= D_{1}(\\mu(D_{2}(x), \\alpha^{s}(y))) - D_{2}(\\mu(D_{1}(x), \\alpha^{k}(y)))\\\\\n&= \\mu(D_{1}(D_{2}(x)), \\alpha^{k + s}(y)) - \\mu(D_{2}(D_{1}(x)), \\alpha^{k + s}(y))\\\\\n&= \\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y)).\n\\end{align*}\nSimilarly, we have $\\nu^{'}(D_{1}, D_{2})(\\mu(x, y)) = \\mu(\\alpha^{k + s}(x), \\nu^{'}(D_{1}, D_{2})(y))$.\n\nHence, we have $\\nu^{'}(D_{1}, D_{2}) \\in C_{\\alpha^{k + s}}(V)$.\n\nTherefore, we have $C(V)$ is a Hom-subalgebras of $(\\mathcal{W}, \\nu^{'}, \\sigma)$.\n\n(2). Suppose that $D_{1} \\in ZDer_{\\alpha^{k}}(V)$, $D_{2} \\in Der_{\\alpha^{s}}(V)$. Then for any $x, y \\in V$,\n\\[\\sigma(D_{1})(\\mu(x, y)) = \\alpha \\circ D_{1}(\\mu(x, y)) = 0.\\]\n\\[\\sigma(D_{1})(\\mu(x, y)) = \\alpha \\circ D_{1}(\\mu(x, y)) = \\alpha(\\mu(D_{1}(x), \\alpha^{k}(y))) = \\mu(\\sigma(D_{1})(x), \\alpha^{k + 1}(y)).\\]\nHence, $\\sigma(D_{1}) \\in ZDer_{\\alpha^{k + 1}}(V)$.\n\\begin{align*}\n&\\nu^{'}(D_{1}, D_{2})(\\mu(x, y))\\\\\n&= D_{1} \\circ D_{2}(\\mu(x, y)) - D_{2} \\circ D_{1}(\\mu(x, y))\\\\\n&= D_{1}(\\mu(D_{2}(x), \\alpha^{s}(y)) + \\mu(\\alpha^{s}(x), D_{2}(y))) = 0.\n\\end{align*}\n\\begin{align*}\n&\\nu^{'}(D_{1}, D_{2})(\\mu(x, y))\\\\\n&= D_{1} \\circ D_{2}(\\mu(x, y)) - D_{2} \\circ D_{1}(\\mu(x, y))\\\\\n&= D_{1}(\\mu(D_{2}(x), \\alpha^{s}(y)) + \\mu(\\alpha^{s}(x), D_{2}(y))) - D_{2}(\\mu(D_{1}(x), \\alpha^{k}(y)))\\\\\n&= \\mu(D_{1}(D_{2}(x)), \\alpha^{k + s}(y)) + \\mu(D_{1}(\\alpha^{s}(x)), \\alpha^{k}(D_{2}(y))) - \\mu(D_{2}(D_{1}(x)), \\alpha^{k + s}(y))\\\\\n&- \\mu(\\alpha^{s}(D_{1}(x)), D_{2}(\\alpha^{k}(y)))\\\\\n&= \\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y)).\n\\end{align*}\nHence, we have $\\nu^{'}(D_{1}, D_{2}) \\in ZDer_{\\alpha^{k + s}}(V)$.\n\nTherefore, $ZDer(V)$ is a Hom-ideal of $Der(V)$.\n\\end{proof}\n\\begin{prop}\\label{prop:3.2}\nLet $(V, \\mu, \\alpha)$ be a multiplicative Hom-Jordan algebra, then\n\\begin{enumerate}[(1)]\n\\item $\\nu^{'}(Der(V), C(V)) \\subseteq C(V)$;\n\\item $\\nu^{'}(QDer(V), QC(V)) \\subseteq QC(V)$;\n\\item $\\nu^{'}(QC(V), QC(V)) \\subseteq QDer(V)$;\n\\item $C(V) \\subseteq QDer(V)$;\n\\item $QDer(V) + QC(V) \\subseteq GDer(V)$;\n\\item $C(V) \\circ Der(V) \\subseteq Der(V)$.\n\\end{enumerate}\n\\end{prop}\n\\begin{proof}\n(1). Suppose that $D_{1} \\in Der_{\\alpha^{k}}(V)$, $D_{2} \\in C_{\\alpha^{s}}(V)$. Then for any $x, y \\in V$, we have\n\\begin{align*}\n&\\nu^{'}(D_{1}, D_{2})(\\mu(x, y))\\\\\n&= D_{1} \\circ D_{2}(\\mu(x, y)) - D_{2} \\circ D_{1}(\\mu(x, y))\\\\\n&= D_{1}(\\mu(D_{2}(x), \\alpha^{s}(y))) - D_{2}(\\mu(D_{1}(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D_{1}(y)))\\\\\n&= \\mu(D_{1}(D_{2}(x)), \\alpha^{k + s}(y)) + \\mu(\\alpha^{k}(D_{2}(x)), D_{1}(\\alpha^{s}(y))) - \\mu(D_{2}(D_{1}(x)), \\alpha^{k + s}(y))\\\\\n&- \\mu(D_{2}(\\alpha^{k}(x)), \\alpha^{s}(D_{1}(y)))\\\\\n&= \\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y)).\n\\end{align*}\nSimilarly, we have $\\nu^{'}(D_{1}, D_{2})(\\mu(x, y)) = \\mu(\\alpha^{k + s}(x), \\nu^{'}(D_{1}, D_{2})(y))$. Hence, $\\nu^{'}(D_{1}, D_{2}) \\in C_{\\alpha^{k + s}}(V)$. Therefore $\\nu^{'}(Der(V), C(V)) \\subseteq C(V)$.\n\n(2). Suppose that $D_{1} \\in QDer_{\\alpha^{k}}(V)$, $D_{2} \\in QC_{\\alpha^{s}}(V)$. Then for any $x, y \\in V$, we have\n\\begin{align*}\n&\\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y))\\\\\n&= \\mu(D_{1} \\circ D_{2}(x), \\alpha^{k + s}(y)) - \\mu(D_{2} \\circ D_{1}(x), \\alpha^{k + s}(y))\\\\\n&= D_{1}^{'}(\\mu(D_{2}(x), \\alpha^{s}(y))) - \\mu(\\alpha^{k}(D_{2}(x)), D_{1}(\\alpha^{s}(y))) - \\mu(\\alpha^{s}(D_{1}(x)), D_{2}(\\alpha^{k}(y)))\\\\\n&= D_{1}^{'}(\\mu(\\alpha^{s}(x), D_{2}(y))) - \\mu(\\alpha^{k}(D_{2}(x)), D_{1}(\\alpha^{s}(y))) - \\mu(\\alpha^{s}(D_{1}(x)), D_{2}(\\alpha^{k}(y)))\\\\\n&= \\mu(D_{1}(\\alpha^{s}(x)), \\alpha^{k}(D_{2}(y))) + \\mu(\\alpha^{k + s}(x), D_{1}(D_{2}(y))) - \\mu(\\alpha^{k}(D_{2}(x)), D_{1}(\\alpha^{s}(y)))\\\\\n&- \\mu(\\alpha^{s}(D_{1}(x)), D_{2}(\\alpha^{k}(y)))\\\\\n&= \\mu(\\alpha^{k + s}(x), D_{1}(D_{2}(y))) - \\mu(\\alpha^{k}(D_{2}(x)), D_{1}(\\alpha^{s}(y)))\\\\\n&= \\mu(\\alpha^{k + s}(x), D_{1}(D_{2}(y))) - \\mu(D_{2}(\\alpha^{k}(x)), \\alpha^{s}(D_{1}(y)))\\\\\n&= \\mu(\\alpha^{k + s}(x), D_{1}(D_{2}(y))) - \\mu(\\alpha^{k + s}(x), D_{2}(D_{1}(y)))\\\\\n&= \\mu(\\alpha^{k + s}(x), \\nu^{'}(D_{1}, D_{2})(y)).\n\\end{align*}\nHence, we have $\\nu^{'}(D_{1}, D_{2}) \\in QC_{\\alpha^{k + s}}(V)$. So $\\nu^{'}(QDer(V), QC(V)) \\subseteq QC(V)$.\n\n(3). Suppose that $D_{1} \\in QC_{\\alpha^{k}}(V)$, $D_{2} \\in QC_{\\alpha^{s}}(V)$. Then for any $x, y \\in V$, we have\n\\begin{align*}\n&\\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y))\\\\\n&= \\mu(D_{1} \\circ D_{2}(x), \\alpha^{k + s}(y)) - \\mu(D_{2} \\circ D_{1}(x), \\alpha^{k + s}(y))\\\\\n&= \\mu(\\alpha^{k}(D_{2}(x)), D_{1}(\\alpha^{s}(y))) - \\mu(\\alpha^{s}(D_{1}(x)), D_{2}(\\alpha^{k}(y)))\\\\\n&= \\mu(D_{2}(\\alpha^{k}(x)), \\alpha^{s}(D_{1}(y))) - \\mu(D_{1}(\\alpha^{s}(x)), \\alpha^{k}(D_{2}(y)))\\\\\n&= \\mu(\\alpha^{k + s}(x), D_{2}(D_{1}(y))) - \\mu(\\alpha^{k + s}(x), D_{1}(D_{2}(y)))\\\\\n&= -\\mu(\\alpha^{k + s}(x), \\nu^{'}(D_{1}, D_{2})(y)),\n\\end{align*}\ni.e., $\\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y)) + \\mu(\\alpha^{k + s}(x), \\nu^{'}(D_{1}, D_{2})(y)) = 0$.\n\nHence, we have $\\nu^{'}(D_{1}, D_{2}) \\in QDer_{\\alpha^{k + s}}(V)$, which implies that $\\nu^{'}(QC(V), QC(V)) \\subseteq QDer(V)$.\n\n(4). Suppose that $D \\in C_{\\alpha^{k}}(V)$. Then for any $x, y \\in V$, we have\n\\[D(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) = \\mu(\\alpha^{k}(x), D(y)).\\]\n\nHence, we have\n\\[\\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D(y)) = 2D(\\mu(x, y)),\\]\nwhich implies that $D \\in QDer_{\\alpha^{k}}(V)$. So $C(V) \\subseteq QDer(V)$.\n\n(5). Suppose that $D_{1} \\in QDer_{\\alpha^{k}}(V)$, $D_{2} \\in QC_{\\alpha^{k}}(V)$. Then for any $x, y \\in V$, we have\n\\begin{align*}\n&\\mu((D_{1} + D_{2})(x), \\alpha^{k}(y))\\\\\n&= \\mu(D_{1}(x), \\alpha^{k}(y)) + \\mu(D_{2}(x), \\alpha^{k}(y))\\\\\n&= D_{1}^{'}(\\mu(x, y)) - \\mu(\\alpha^{k}(x), D_{1}(y)) + \\mu(\\alpha^{k}(x), D_{2}(y))\\\\\n&= D_{1}^{'}(\\mu(x, y)) - \\mu(\\alpha^{k}(x), (D_{1} - D_{2})(y)).\n\\end{align*}\nSince $D_{1}^{'}$, $D_{1} - D_{2} \\in End(V)$, we have $D_{1} + D_{2} \\in GDer_{\\alpha^{k}}(V)$. So $QDer(V) + QC(V) \\subseteq GDer(V)$.\n\n(6). Suppose that $D_{1} \\in C_{\\alpha^{k}}(V)$, $D_{2} \\in Der_{\\alpha^{s}}(V)$. Then for any $x, y \\in V$, we have\n\\begin{align*}\n&D_{1} \\circ D_{2}(\\mu(x, y))\\\\\n&= D_{1}(\\mu(D_{2}(x), \\alpha^{s}(y)) + \\mu(\\alpha^{s}(x), D_{2}(y)))\\\\\n&= \\mu(D_{1}(D_{2}(x)), \\alpha^{k + s}(y)) + \\mu(\\alpha^{k + s}(x), D_{1}(D_{2}(y))),\n\\end{align*}\nwhich implies that $D_{1} \\circ D_{2} \\in Der_{\\alpha^{k + s}}(V)$. So $C(V) \\circ Der(V) \\subseteq Der(V)$.\n\\end{proof}\n\\begin{thm}\nSuppose that $(V, \\mu, \\alpha)$ is a multiplicative Hom-Jordan algebra. Then $GDer(V) = QDer(V) + QC(V)$.\n\\end{thm}\n\\begin{proof}\nAccording to Proposition \\ref{prop:3.2} (5), we need only to show that $GDer(V) \\subseteq QDer(V) + QC(V)$.\n\nSuppose that $D \\in GDer_{\\alpha^{k}}(V)$. Then there exist $D^{'},\\; D^{''} \\in End(V)$ such that\n\\[D^{''}(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D^{'}(y)),\\quad\\forall x, y \\in V.\\]\n\nSince $\\mu$ is commutative, we have\n\\[D^{''}(\\mu(y, x)) = \\mu(D^{'}(y), \\alpha^{k}(x)) + \\mu(\\alpha^{k}(y), D(x)),\\quad\\forall x, y \\in V,\\]\nwhich implies that $D^{'} \\in GDer_{\\alpha^{k}}(V)$.\n\nMoreover, we have\n\\begin{align*}\n&\\mu\\left(\\frac{D + D^{'}}{2}(x), \\alpha^{k}(y)\\right) + \\mu\\left(\\alpha^{k}(x), \\frac{D + D^{'}}{2}(y)\\right)\\\\\n&= \\frac{1}{2}(\\mu(D(x), \\alpha^{k}(y)) + \\mu(D^{'}(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D(y)) + \\mu(\\alpha^{k}(x), D^{'}(y)))\\\\\n&= D^{''}(\\mu(x, y)),\n\\end{align*}\nwhich implies $\\frac{D + D^{'}}{2} \\in QDer_{\\alpha^{k}}(V)$.\n\\begin{align*}\n&\\mu\\left(\\frac{D - D^{'}}{2}(x), \\alpha^{k}(y)\\right)\\\\\n&= \\frac{1}{2}(\\mu(D(x), \\alpha^{k}(y)) - \\mu(D^{'}(x), \\alpha^{k}(y)))\\\\\n&= \\frac{1}{2}(D^{''}(\\mu(x, y)) - \\mu(\\alpha^{k}(x), D^{'}(y)) - D^{''}(\\mu(x, y)) + \\mu(\\alpha^{k}(x), D(y)))\\\\\n&= \\mu\\left(\\alpha^{k}(x), \\frac{D - D^{'}}{2}(y)\\right),\n\\end{align*}\nwhich implies that $\\frac{D - D^{'}}{2} \\in QC_{\\alpha^{k}}(V)$.\n\nHence, $D = \\frac{D + D^{'}}{2} + \\frac{D - D^{'}}{2} \\in QDer(V) + QC(V)$, i.e., $GDer(V) \\subseteq QDer(V) + QC(V)$.\n\nTherefore, $GDer(V) = QDer(V) + QC(V)$.\n\\end{proof}\n\\begin{prop}\nSuppose that $(V, \\mu, \\alpha)$ is a multiplicative Hom-Jordan algebra where $V$ can be decomposed into the direct sum of two Hom-ideals, i.e., $V = V_{1} \\oplus V_{2}$ and $\\alpha$ is surjective. Then we have\n\\begin{enumerate}[(1)]\n\\item $Z(V) = Z(V_{1}) \\oplus Z(V_{2})$;\n\\item If\\; $Z(V) = \\{0\\}$, then we have\n \\begin{enumerate}[(a)]\n \\item $Der(V) = Der(V_{1}) \\oplus Der(V_{2})$;\n \\item $GDer(V) = GDer(V_{1}) \\oplus GDer(V_{2})$;\n \\item $QDer(V) = QDer(V_{1}) \\oplus QDer(V_{2})$;\n \\item $C(V) = C(V_{1}) \\oplus C(V_{2})$.\n \\end{enumerate}\n\\end{enumerate}\n\\end{prop}\n\\begin{proof}\n(1). Obviously, $Z(V_{1}) \\cap Z(V_{2}) = \\{0\\}$. And it's easy to show that $Z(V_{i})(i = 1, 2)$ are Hom-ideals of $Z(V)$. For all $z_{i} \\in Z(V_{i})(i = 1, 2)$, take $x = x_{1} + x_{2}$ where $x_{1} \\in J_{1}$, $x_{2} \\in J_{2}$. We have\n\\[\\mu(z_{1} + z_{2}, x) = \\mu(z_{1} + z_{2}, x_{1} + x_{2}) = \\mu(z_{1}, x_{1}) + \\mu(z_{2}, x_{2}),\\]\nsince $z_{i} \\in Z(V_{i})(i = 1, 2)$, we have\n\\[\\mu(z_{i}, x_{i}) = 0,\\; i = 1, 2.\\]\nHence,\n\\[\\mu(z_{1} + z_{2}, x) = 0,\\]\nwhich implies that $z_{1} + z_{2} \\in Z(V)$, i.e., $Z(V_{1}) \\oplus Z(V_{2}) \\subseteq Z(V)$.\n\nOn the other hand, for all $a \\in Z(V)$, suppose that $a = a_{1} + a_{2}$ where $a_{i} \\in V_{i}(i = 1, 2)$. Then for all $x_{1} \\in V_{1}$,\n\\[\\mu(a_{1}, x_{1}) = \\mu(a - a_{2}, x_{1}) = \\mu(a, x_{1})\\]\nsince $a \\in Z(V)$, we have\n\\[\\mu(a, x_{1}) = 0.\\]\nHence,\n\\[\\mu(a_{1}, x_{1}) = 0,\\]\nwhich implies that $a_{1} \\in Z(V_{1})$. Similarly, we have $a_{2} \\in Z(V_{2})$. Hence, $Z(V) \\subseteq Z(V_{1}) \\oplus Z(V_{2})$.\n\nTherefore, we have $Z(V) = Z(V_{1}) \\oplus Z(V_{2})$.\n\n(2). $\\bf{Step 1}$. We'll show that $\\forall i = 1, 2$, $D(V_{i}) \\subseteq V_{i},\\quad\\forall D \\in Der_{\\alpha^{k}}(V)(k \\geq 0)$.\n\nSuppose that $x_{i} \\in V_{i}$, then\n\\[\\mu(D(x_{1}), \\alpha^{k}(x_{2})) = D(\\mu(x_{1}, x_{2})) - \\mu(\\alpha^{k}(x_{1}), D(x_{2})) \\in V_{1} \\cap V_{2} = 0,\\]\nsince $V_{1}$, $V_{2}$ are Hom-ideals of $(V, \\mu, \\alpha)$.\n\nSuppose that $D(x_{1}) = u_{1} + u_{2}$ where $u_{1} \\in V_{1}$, $u_{2} \\in V_{2}$. Then\n\\[\\mu(u_{2}, \\alpha^{k}(x_{2})) = \\mu(u_{1} + u_{2}, \\alpha^{k}(x_{2})) = \\mu(D(x_{1}), \\alpha^{k}(x_{2})) = 0,\\]\nwhich implies that $u_{2} \\in Z(V_{2})$ since $\\alpha$ is surjective. Note that $Z(V) = \\{0\\}$, we have $Z(V_{i}) = \\{0\\}(i = 1, 2)$. Hence, $u_{2} = 0$. That is to say $D(x_{1}) \\in V_{1}$. Similarly, we have $D(x_{2}) \\in V_{2}$.\n\n$\\bf{Step 2}$. We'll show that $Der_{\\alpha^{k}}(V_{1}) \\dotplus Der_{\\alpha^{k}}(V_{2}) \\subseteq Der_{\\alpha^{k}}(V)(k \\geq 0)$.\n\nFor any $D \\in Der_{\\alpha^{k}}(V_{1})$, we extend it to a linear map on $V$ as follow\n\\[D(x_{1} + x_{2}) = D(x_{1}),\\quad\\forall x_{1} \\in V_{1}, x_{2} \\in V_{2}.\\]\n\nThen for any $x, y \\in V$, suppose that $x = x_{1} + x_{2}$, $y = y_{1} + y_{2} \\in V$, where $x_{1}, y_{1} \\in V_{1}$, $x_{2}, y_{2} \\in V_{2}$, we have\n\\[D(\\mu(x, y)) = D(\\mu(x_{1} + x_{2}, y_{1} + y_{2})) = D(\\mu(x_{1}, y_{1}) + \\mu(x_{2}, y_{2})) = D(\\mu(x_{1}, y_{1})),\\]\n\\begin{align*}\n&\\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D(y))\\\\\n&= \\mu(D(x_{1} + x_{2}), \\alpha^{k}(y_{1} + y_{2})) + \\mu(\\alpha^{k}(x_{1} + x_{2}), D(y_{1} + y_{2}))\\\\\n&= \\mu(D(x_{1}), \\alpha^{k}(y_{1})) + \\mu(\\alpha^{k}(x_{1}), D(y_{1})),\n\\end{align*}\nsince $D \\in Der_{\\alpha^{k}}(V_{1})$,\n\\[D(\\mu(x_{1}, y_{1})) = \\mu(D(x_{1}), \\alpha^{k}(y_{1})) + \\mu(\\alpha^{k}(x_{1}), D(y_{1})),\\]\nwe have\n\\[D(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D(y)),\\]\nwhich implies that $D \\in Der_{\\alpha^{k}}(V)$, i.e., $Der_{\\alpha^{k}}(V_{1}) \\subseteq Der_{\\alpha^{k}}(V)$. Moreover, $D \\in Der_{\\alpha^{k}}(V_{1})$ if and only if $D(x_{2}) = 0,\\; \\forall x_{2} \\in V_{2}$.\n\nSimilarly, we have $Der_{\\alpha^{k}}(V_{2}) \\subseteq Der_{\\alpha^{k}}(V)$ and $D \\in Der_{\\alpha^{k}}(V_{2})$ if and only if $D(x_{1}) = 0,\\; \\forall x_{1} \\in V_{1}$.\n\nThen we have $Der_{\\alpha^{k}}(V_{1}) + Der_{\\alpha^{k}}(V_{2}) \\subseteq Der_{\\alpha^{k}}(V)$ and $Der_{\\alpha^{k}}(V_{1}) \\cap Der_{\\alpha^{k}}(V_{2}) = \\{0\\}$. Hence, $Der_{\\alpha^{k}}(V_{1}) \\dotplus Der_{\\alpha^{k}}(V_{2}) \\subseteq Der_{\\alpha^{k}}(V)$.\n\n$\\bf{Step 3}$. We'll prove that $Der_{\\alpha^{k}}(V_{1}) \\dotplus Der_{\\alpha^{k}}(V_{2}) = Der_{\\alpha^{k}}(V)$.\n\nSuppose that $D \\in Der_{\\alpha^{k}}(V)$. Set $x = x_{1} + x_{2}, x_{i} \\in V_{i}$. Define $D_{1}, D_{2}$ as follows\n\\begin{equation}\n\\left\\{\n\\begin{aligned}\nD_{1}(x_{1} + x_{2}) = D(x_{1}),\\\\\nD_{2}(x_{1} + x_{2}) = D(x_{2}).\n\\end{aligned}\n\\right.\n\\end{equation}\nObviously, $D = D_{1} + D_{2}$.\n\nFor any $u_{1}, v_{1} \\in V_{1}$,\n\\begin{align*}\n&D_{1}(\\mu(u_{1}, v_{1})) = D(\\mu(u_{1}, v_{1})) = \\mu(D(u_{1}), \\alpha^{k}(v_{1})) + \\mu(\\alpha^{k}(u_{1}), D(v_{1}))\\\\\n&= \\mu(D_{1}(u_{1}), \\alpha^{k}(v_{1})) + \\mu(\\alpha^{k}(u_{1}), D_{1}(v_{1})).\n\\end{align*}\nHence, $D_{1} \\in Der_{\\alpha^{k}}(V_{1})$. Similarly, $D_{2} \\in Der_{\\alpha^{k}}(V_{2})$.\n\nTherefore, $Der_{\\alpha^{k}}(V_{1}) \\dotplus Der_{\\alpha^{k}}(V_{2}) = Der_{\\alpha^{k}}(V)$ as a vector space.\n\nHence, we have\n\\begin{align*}\n&Der(V) = \\dotplus_{k \\geq 0}Der_{\\alpha^{k}}(V) = \\dotplus_{k \\geq 0}(Der_{\\alpha^{k}}(V_{1}) \\dotplus Der_{\\alpha^{k}}(V_{2}))\\\\\n&= (\\dotplus_{k \\geq 0}Der_{\\alpha^{k}}(V_{1})) \\dotplus (\\dotplus_{k \\geq 0}Der_{\\alpha^{k}}(V_{2})) = Der(V_{1}) \\dotplus Der(V_{2}).\n\\end{align*}\n\n$\\bf{Step 4}$. We'll show that $Der(V_{i})(i = 1, 2)$ are Hom-ideals of $Der(V)$.\nSuppose that $D_{1} \\in Der_{\\alpha^{k}}(V_{1})$, $D_{2} \\in Der_{\\alpha^{s}}(V)$. Then for any $x_{1}, y_{1} \\in V_{1}$, we have\n\\begin{align*}\n&\\sigma(D_{1})(\\mu(x_{1}, y_{1})) = \\alpha \\circ D_{1}(\\mu(x_{1}, y_{1})) = \\alpha(\\mu(D_{1}(x_{1}), \\alpha^{k}(y_{1})) + \\mu(\\alpha^{k}(x_{1}), D_{1}(y_{1})))\\\\\n&= \\mu(\\sigma(D_{1})(x_{1}), \\alpha^{k + 1}(y_{1})) + \\mu(\\alpha^{k + 1}(x_{1}), \\sigma(D_{1})(y_{1})),\n\\end{align*}\nwhich implies that $\\sigma(D_{1}) \\in Der_{\\alpha^{k + 1}}(V_{1})$.\n\\begin{align*}\n&\\nu^{'}(D_{1}, D_{2})(\\mu(x_{1}, y_{1}))\\\\\n&= D_{1} \\circ D_{2}(\\mu(x_{1}, y_{1})) - D_{2} \\circ D_{1}(\\mu(x_{1}, y_{1}))\\\\\n&= D_{1}(\\mu(D_{2}(x_{1}), \\alpha^{s}(y_{1})) + \\mu(\\alpha^{s}(x_{1}), D_{2}(y_{1}))) - D_{2}(\\mu(D_{1}(x_{1}), \\alpha^{k}(y_{1})) + \\mu(\\alpha^{k}(x_{1}), D_{1}(y_{1})))\\\\\n&= \\mu(D_{1}(D_{2}(x_{1})), \\alpha^{k + s}(y_{1})) + \\mu(\\alpha^{k}(D_{2}(x_{1})), D_{1}(\\alpha^{s}(y_{1}))) + \\mu(D_{1}(\\alpha^{s}(x_{1})), \\alpha^{k}(D_{2}(y_{1})))\\\\\n&+ \\mu(\\alpha^{k + s}(x_{1}), D_{1}(D_{2}(y_{1}))) - \\mu(D_{2}(D_{1}(x_{1})), \\alpha^{k + s}(y_{1})) - \\mu(\\alpha^{s}(D_{1}(x_{1})), D_{2}(\\alpha^{k}(y_{1})))\\\\\n&- \\mu(D_{2}(\\alpha^{k}(x_{1})), \\alpha^{s}(D_{1}(y_{1}))) - \\mu(\\alpha^{k + s}(x_{1}), D_{2}(D_{1}(y_{1})))\\\\\n&= \\mu(\\nu^{'}(D_{1}, D_{2})(x_{1}), \\alpha^{k + s}(y_{1})) + \\mu(\\alpha^{k + s}(x_{1}), \\nu^{'}(D_{1}, D_{2})(y_{1})),\n\\end{align*}\nwhich implies that $\\nu^{'}(D_{1}, D_{2}) \\in Der_{\\alpha^{k + s}}(V_{1})$. So $Der(V_{1})$ is a Hom-ideal of $Der(V)$.\n\nSimilarly, we have $Der(V_{2})$ is a Hom-ideal of $Der(V)$.\n\nTherefore, $Der(V) = Der(V_{1}) \\oplus Der(V_{2})$.\n\n(3), (4), (5) similar to the proof of (2).\n\\end{proof}\n\\begin{thm}\nLet $(V, \\mu, \\alpha)$ be a multiplicative Hom-Jordan algebra, $\\alpha$ a surjection and $Z(V)$ the centralizer of $(V, \\mu, \\alpha)$. Then $\\nu^{'}(C(V), QC(V)) \\subseteq End(V, Z(V))$. Moreover, if $Z(V) = \\{0\\}$, then $\\nu^{'}(C(V), QC(V)) = \\{0\\}$.\n\\end{thm}\n\\begin{proof}\nSuppose that $D_{1} \\in C_{\\alpha^{k}}(V)$, $D_{2} \\in QC_{\\alpha^{s}}(V)$ and $x \\in V$. Since $\\alpha$ is surjective, there exists $y^{'} \\in V$ such that $y = \\alpha^{k + s}(y^{'})$ for any $y \\in V$.\n\\begin{align*}\n&\\mu(\\nu^{'}(D_{1}, D_{2})(x), y)\\\\\n&= \\mu(D_{1} \\circ D_{2}(x), \\alpha^{k + s}(y^{'})) - \\mu(D_{2} \\circ D_{1}(x), \\alpha^{k + s}(y^{'}))\\\\\n&= D_{1}(\\mu(D_{2}(x), \\alpha^{s}(y^{'}))) - \\mu(\\alpha^{s}(D_{1}(x)), D_{2}(\\alpha^{k}(y^{'})))\\\\\n&= D_{1}(\\mu(D_{2}(x), \\alpha^{s}(y^{'}))) - \\mu(D_{1}(\\alpha^{s}(x)), \\alpha^{k}(D_{2}(y^{'})))\\\\\n&= D_{1}(\\mu(D_{2}(x), \\alpha^{s}(y^{'}))) - D_{1}(\\mu(\\alpha^{s}(x), D_{2}(y^{'})))\\\\\n&= D_{1}(\\mu(D_{2}(x), \\alpha^{s}(y^{'})) - \\mu(\\alpha^{s}(x), D_{2}(y^{'})))\\\\\n&= 0,\n\\end{align*}\nwhich implies that $\\nu^{'}(D_{1}, D_{2})(x) \\in Z(V)$, i.e., $\\nu^{'}(C(V), QC(V)) \\subseteq End(V, Z(V))$.\n\nFurthermore, if $Z(V) = \\{0\\}$, it's clearly that $\\nu^{'}(C(V), QC(V)) = \\{0\\}$.\n\\end{proof}\n\\begin{thm}\nLet $(V, \\mu, \\alpha)$ be a multiplicative Hom-Jordan algebra. Then $ZDer(V) = C(V) \\cap Der(V)$.\n\\end{thm}\n\\begin{proof}\nAssume that $D \\in C(V)_{\\alpha^{k}} \\cap Der_{\\alpha^{k}}(V)$. Then for any $x, y \\in V$, we have\n\\[D(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D(y))\\]\nsince $D \\in Der_{\\alpha^{k}}(V)$.\n\\[D(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) = \\mu(\\alpha^{k}(x), D(y))\\]\nsince $D \\in C_{\\alpha^{k}}(V)$.\n\nHence we have $D(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) = 0$, which implies that $D \\in ZDer_{\\alpha^{k}}(V)$. Therefore, $C(V) \\cap Der(V) \\subseteq ZDer(V)$.\n\nOn the other hand, assume that $D \\in ZDer_{\\alpha^{k}}(V)$. Then for any $x, y \\in V$, we have\n\\[D(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) = 0,\\]\nhence we have\n\\[\\mu(\\alpha^{k}(x), D(y)) = \\mu(D(y), \\alpha^{k}(x)) = 0.\\]\nTherefore\n\\[D(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D(y))\\]\nwhich implies that $D \\in Der_{\\alpha^{k}}(V)$.\n\nAnd\n\\[D(\\mu(x, y)) = \\mu(D(x), \\alpha^{k}(y)) = \\mu(\\alpha^{k}(x), D(y))\\]\nwhich implies that $D \\in C_{\\alpha^{k}}(V)$.\n\nTherefore $D \\in C(V)_{\\alpha^{k}} \\cap Der_{\\alpha^{k}}(V)$, i.e., $ZDer(V) \\subseteq C(V) \\cap Der(V)$.\n\nHence, we have $ZDer(V) = C(V) \\cap Der(V)$.\n\\end{proof}\n\\begin{thm}\\label{thm:3.7}\nLet $(V, \\mu, \\alpha)$ be a multiplicative Hom-Jordan algebra. Then $(QC(V), \\nu, \\sigma)$ is a Hom-Jordan algebra.\n\\end{thm}\n\\begin{proof}\nAccording to Lemma \\ref{le:2.7} (1), we need only to show that $QC(V)$ is a Hom-subalgebra of $(\\mathcal{W}, \\nu, \\sigma)$.\n\nSuppose that $D_{1} \\in QC_{\\alpha^{k}}(V)$, $D_{2} \\in QC_{\\alpha^{s}}(V)$. Then for any $x, y \\in V$, we have\n\\begin{align*}\n&\\mu(\\sigma(D_{1})(x), \\alpha^{k + 1}(y)) = \\mu(\\alpha \\circ D_{1}(x), \\alpha^{k + 1}(y)) = \\alpha(\\mu(D_{1}(x), \\alpha^{k}(y)))\\\\\n&= \\alpha(\\mu(\\alpha^{k}(x), D_{1}(y))) = \\mu(\\alpha^{k + 1}(x), \\sigma(D_{1})(y)),\n\\end{align*}\nwhich implies that $\\sigma(D_{1}) \\in QC_{\\alpha^{k + 1}}(V)$.\n\\begin{align*}\n&\\mu(\\nu(D_{1}, D_{2})(x), \\alpha^{k + s}(y))\\\\\n&= \\mu(D_{1} \\circ D_{2}(x), \\alpha^{k + s}(y)) + \\mu(D_{2} \\circ D_{1}(x), \\alpha^{k + s}(y))\\\\\n&= \\mu(\\alpha^{k}(D_{2}(x)), D_{1}(\\alpha^{s}(y))) + \\mu(\\alpha^{s}(D_{1}(x)), D_{2}(\\alpha^{k}(y)))\\\\\n&= \\mu(D_{2}(\\alpha^{k}(x)), \\alpha^{s}(D_{1}(y))) + \\mu(D_{1}(\\alpha^{s}(x)), \\alpha^{k}(D_{2}(y)))\\\\\n&= \\mu(\\alpha^{k + s}(x), D_{2}(D_{1}(y))) + \\mu(\\alpha^{k + s}(x), D_{1}(D_{2}(y)))\\\\\n&= \\mu(\\alpha^{k + s}(x), \\nu(D_{1}, D_{2})(y)),\n\\end{align*}\nwhich implies that $\\nu(D_{1}, D_{2}) \\in QC_{\\alpha^{k + s}}(V)$. Therefore, $QC(V)$ is a Hom-subalgebra of $(\\mathcal{W}, \\nu, \\sigma)$, i.e., $(QC(V), \\nu, \\sigma)$ is a Hom-Jordan algebra.\n\\end{proof}\n\\begin{thm}\nSuppose that $(V, \\mu, \\alpha)$ be a multiplicative Hom-Jordan algebra over a field $\\rm{F}$ of characteristic other than $2$.\n\\begin{enumerate}[(1)]\n\\item $(QC(V), \\nu^{'}, \\sigma)$ is a Hom-Lie algebra if and only if $(QC(V), \\iota, \\sigma)$ is a Hom-associative algebra where $\\iota(D_{1}, D_{2}) = D_{1} \\circ D_{2},\\quad\\forall D_{1}, D_{2} \\in QC(V)$.\n\\item If $\\alpha$ is a surjection and $Z(V) = \\{0\\}$, then $(QC(V), \\nu^{'}, \\sigma)$ is a Hom-Lie algebra if and only if $\\nu^{'}(QC(V), QC(V)) = \\{0\\}$.\n\\end{enumerate}\n\\end{thm}\n\\begin{proof}\n(1). $(\\Leftarrow)$. For all $D_{1}, D_{2} \\in QC(V)$, we have $D_{1} \\circ D_{2} \\in QC(V)$, $D_{2} \\circ D_{1} \\in QC(V)$. So $\\nu^{'}(D_{1}, D_{2}) \\in QC(V)$. Moreover, $\\sigma(D_{1}) \\in QC(V)$. Therefore, $(QC(V), \\nu^{'}, \\sigma)$ is a Hom-Lie algebra.\n\n$(\\Rightarrow)$. For all $D_{1}, D_{2} \\in QC(V)$, we have $\\iota(D_{1}, D_{2}) = \\frac{1}{2}(\\nu^{'}(D_{1}, D_{2}) + \\nu(D_{1}, D_{2}))$. According to Theorem \\ref{thm:3.7}, we have $\\nu(D_{1}, D_{2}) \\in QC(V)$. Note that $\\nu^{'}(D_{1}, D_{2}) \\in QC(V)$, we have $\\iota(D_{1}, D_{2}) \\in QC(V)$. Hence, $(QC(V), \\iota, \\sigma)$ is a Hom-associative algebra.\n\n(2). $(\\Rightarrow)$. Suppose that $D_{1} \\in QC_{\\alpha^{k}}(V)$, $D_{2} \\in QC_{\\alpha^{s}}(V)$ and $x \\in V$. Since $\\alpha$ is a surjection, there exists $y^{'} \\in V$ such that $y = \\alpha^{k + s}(y^{'})$ for any $y \\in V$. Note that $(QC(V), \\nu^{'}, \\sigma)$ is a Hom-Lie algebra, then $\\nu^{'}(D_{1}, D_{2}) \\in QC_{\\alpha^{k + s}}(V)$. Then\n\\[\\mu(\\nu^{'}(D_{1}, D_{2})(x), y) = \\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y^{'})) = \\mu(\\alpha^{k + s}(x), \\nu^{'}(D_{1}, D_{2})(y^{'})).\\]\n\nNote that\n\\begin{align*}\n&\\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y^{'}))\\\\\n&= \\mu(D_{1} \\circ D_{2}(x), \\alpha^{k + s}(y^{'})) - \\mu(D_{2} \\circ D_{1}(x), \\alpha^{k + s}(y^{'}))\\\\\n&= \\mu(\\alpha^{k}(D_{2}(x)), D_{1}(\\alpha^{s}(y^{'}))) - \\mu(\\alpha^{s}(D_{1}(x)), D_{2}(\\alpha^{k}(y^{'})))\\\\\n&= \\mu(D_{2}(\\alpha^{k}(x)), \\alpha^{s}(D_{1}(y^{'}))) - \\mu(D_{1}(\\alpha^{s}(x)), \\alpha^{k}(D_{2}(y^{'})))\\\\\n&= \\mu(\\alpha^{k + s}(x), D_{2} \\circ D_{1}(y^{'})) - \\mu(\\alpha^{k + s}(x), D_{1} \\circ D_{2}(y^{'}))\\\\\n&= -\\mu(\\alpha^{k + s}(x), \\nu^{'}(D_{1}, D_{2})(y^{'})).\n\\end{align*}\nHence, we have $\\mu(\\nu^{'}(D_{1}, D_{2})(x), y) = \\mu(\\nu^{'}(D_{1}, D_{2})(x), \\alpha^{k + s}(y^{'})) = 0$, which implies that $\\nu^{'}(D_{1}, D_{2})(x) \\in Z(V)$. Note that $Z(V) = \\{0\\}$, we have $\\nu^{'}(D_{1}, D_{2})(x) = 0$, i.e., $\\nu^{'}(D_{1}, D_{2}) = 0$. Therefore, $\\nu^{'}(QC(V), QC(V)) = \\{0\\}$.\n\n$(\\Leftarrow)$. Obviously.\n\\end{proof}\n\n\\section{The quasiderivations of Hom-Jordan algebras}\\label{se:4}\n\nIn this section, we will prove that the quasiderivations of $(V, \\mu, \\alpha)$ can be embedded as derivations in a larger Hom-Jordan agebra and obtain a direct sum decomposition of $Der(\\breve{V})$ when the centralizer of $(V, \\mu, \\alpha)$ is equal to zero.\n\\begin{prop}\nSuppose that $(V, \\mu, \\alpha)$ is a Hom-Jordan algebra over $\\rm{F}$ and $t$ is an indeterminate. We define $\\breve{V} := \\{\\sum(x \\otimes t + y \\otimes t^{2}) | x, y \\in V\\}$, $\\breve{\\alpha}(\\breve{V}) := \\{\\sum(\\alpha(x) \\otimes t + \\alpha(y) \\otimes t^{2}) | x, y \\in V\\}$ and $\\breve{\\mu}(x \\otimes t^{i}, y \\otimes t^{j}):= \\mu(x,y) \\otimes t^{i + j}$ where $x, y \\in V, i, j \\in \\{1, 2\\}$. Then $(\\breve{V}, \\breve{\\mu}, \\breve{\\alpha})$ is a Hom-Jordan algebra.\n\\end{prop}\n\\begin{proof}\nIt's obvious that $\\breve{\\mu}$ is a bilinear map and commutative since $\\mu$ is a bilinear map and commutative.\n\nFor any $x \\otimes t^{i}, y \\otimes t^{j} \\in \\breve{V}$, we have\n\\begin{align*}\n&\\breve{\\mu}(\\breve{\\alpha}^{2}(x \\otimes t^{i}), \\breve{\\mu}(y \\otimes t^{j}, \\breve{\\mu}(x \\otimes t^{i}, x \\otimes t^{i})))\\\\\n&= \\mu(\\alpha^{2}(x), \\mu(y, \\mu(x, x))) \\otimes t^{3i + j}\\\\\n&= \\mu(\\mu(\\alpha(x), y), \\alpha(\\mu(x, x))) \\otimes t^{3i + j}\\\\\n&= \\breve{\\mu}(\\breve{\\mu}(\\breve{\\alpha}(x \\otimes t^{i}), y \\otimes t^{j}), \\breve{\\alpha}(\\breve{\\mu}(x \\otimes t^{i}, x \\otimes t^{i}))).\n\\end{align*}\nTherefore, $(\\breve{V}, \\breve{\\mu}, \\breve{\\alpha})$ is a Hom-Jordan algebra.\n\\end{proof}\n\nFor convenience, we write $xt(xt^{2})$ in place of $x \\otimes t(x \\otimes t^{2})$.\n\nIf $U$ is a subspace of $V$ such that $V = U \\dotplus \\mu(V, V)$, then\n\\[\\breve{V} = Vt + Vt^{2} = Vt + Ut^{2} + \\mu(V, V)t^{2}.\\]\nNow we define a map $\\varphi : QDer(V) \\rightarrow End(\\breve{V})$ satisfying\n\\[\\varphi(D)(at + ut^{2} + bt^{2}) = D(a)t + D^{'}(b)t^{2},\\]\nwhere $D \\in QDer_{\\alpha^{k}}(V)$, and $D^{'}$ is in Definition \\ref{defn:2.11} (2), $a \\in V, u \\in U, b \\in \\mu(V, V)$.\n\\begin{prop}\\label{prop:4.2}\n$V, \\breve{V}, \\varphi$ are defined as above. Then\n\\begin{enumerate}[(1)]\n\\item $\\varphi$ is injective and $\\varphi(D)$ does not depend on the choice of $D^{'}$.\n\\item $\\varphi(QDer(V)) \\subseteq Der(\\breve{V})$.\n\\end{enumerate}\n\\end{prop}\n\\begin{proof}\n(1). If $\\varphi(D_{1}) = \\varphi(D_{2})$, then for all $a \\in V, u \\in U, b \\in \\mu(V, V)$, we have\n\\[\\varphi(D_{1})(at + ut^{2} + bt^{2}) = \\varphi(D_{2})(at + ut^{2} + bt^{2}),\\]\nwhich implies that\n\\[D_{1}(a)t + D_{1}^{'}(b)t^{2} = D_{2}(a)t + D_{2}^{'}(b)t^{2}.\\]\nHence we have\n\\[D_{1}(a) = D_{2}(a),\\]\nwhich implies that $D_{1} = D_{2}$. Therefore, $\\varphi$ is injective.\n\nSuppose that there exists $D^{''}$ such that $\\varphi(D)(at + ut^{2} + bt^{2}) = D(a)t + D^{''}(b)t^{2}$ and $\\mu(D(x), \\alpha^{k}(y)) + \\mu(\\alpha^{k}(x), D(y)) = D^{''}(\\mu(x, y))$, then we have\n\\[D^{''}(\\mu(x, y)) = D^{'}(\\mu(x, y)),\\]\nwhich implies that $D^{''}(b) = D^{'}(b)$.\n\nHence, we have\n\\[\\varphi(D)(at + ut^{2} + bt^{2}) = D(a)t + D^{''}(b)t^{2} = D(a)t + D^{'}(b)t^{2}.\\]\nThat is to say $\\varphi(D)$ does not depend on the choice of $D^{'}$.\n\n(2). $\\forall i + j \\geq 3$, we have $\\breve{\\mu}(x \\otimes t^{i}, y \\otimes t^{j}) = \\mu(x, y) \\otimes t^{i + j} = 0$. Hence, we need only to show that\n\\[\\varphi(D)(\\breve{\\mu}(xt, yt)) = \\breve{\\mu}(\\varphi(D)(xt), \\breve{\\alpha}^{k}(yt)) + \\breve{\\mu}(\\breve{\\alpha}^{k}(xt), \\varphi(D)(yt)).\\]\n\nFor all $x, y \\in V$, we have\n\\begin{align*}\n&\\varphi(D)(\\breve{\\mu}(xt, yt))\\\\\n&= \\varphi(D)(\\mu(x, y)t^{2}) = D^{'}(\\mu(x, y))t^{2}\\\\\n&= \\mu(D(x), \\alpha^{k}(y))t^{2} + \\mu(\\alpha^{k}(x), D(y))t^{2}\\\\\n&= \\breve{\\mu}(D(x)t, \\alpha^{k}(y)t) + \\breve{\\mu}(\\alpha^{k}(x)t, D(y)t)\\\\\n&= \\breve{\\mu}(\\varphi(D)(xt), \\breve{\\alpha}^{k}(yt)) + \\breve{\\mu}(\\breve{\\alpha}^{k}(xt), \\varphi(D)(yt)).\n\\end{align*}\nHence, $\\varphi(D) \\in Der_{\\breve{\\alpha}^{k}}(\\breve{V})$. Therefore, $\\varphi(QDer(V)) \\subseteq Der(\\breve{V})$.\n\\end{proof}\n\\begin{prop}\nLet $(V, \\mu, \\alpha)$ be a Hom-Jordan algebra. $Z(V) = \\{0\\}$ and $\\breve{V}, \\varphi$ are defined as above. Then $Der(\\breve{V}) = \\varphi(QDer(V)) \\dotplus ZDer(\\breve{V})$.\n\\end{prop}\n\\begin{proof}\nAssume that $xt + yt^{2} \\in Z(\\breve{V})$, then for any $x^{'}t + y^{'}t^{2} \\in \\breve{V}$, we have\n\\[0 = \\breve{\\mu}(xt + yt^{2}, x^{'}t + y^{'}t^{2}) = \\mu(x, x^{'})t^{2},\\]\nwhich implies that $\\mu(x, x^{'}) = 0$. Note that $Z(V) = \\{0\\}$, we have $x = 0$. Hence, $Z(\\breve{V}) \\subseteq Vt^{2}$. Obviously, $Vt^{2} \\subseteq Z(\\breve{V})$. Therefore, $Z(\\breve{V}) = Vt^{2}$.\n\nSuppose that $g \\in Der_{\\breve{\\alpha}^{k}}(\\breve{V})$, $at^{2} \\in Z(\\breve{V})$. Since $\\alpha$ is surjective, $\\breve{\\alpha}$ is also surjective. For any $xt + yt^{2} \\in \\breve{V}$, there exists $x^{'}t + y^{'}t^{2} \\in \\breve{V}$ such that $xt + yt^{2} = \\breve{\\alpha}^{k}(x^{'}t + y^{'}t^{2})$.\n\\begin{align*}\n&\\breve{\\mu}(g(at^{2}), xt + yt^{2}) = \\breve{\\mu}(g(at^{2}), \\breve{\\alpha}^{k}(x^{'}t + y^{'}t^{2}))\\\\\n&= g(\\breve{\\mu}(at^{2}, x^{'}t + y^{'}t^{2})) - \\breve{\\mu}(\\breve{\\alpha}^{k}(at^{2}), g(x^{'}t + y^{'}t^{2}))\\\\\n&= 0,\n\\end{align*}\nwhich implies that $g(at^{2}) \\in Z(\\breve{V})$. Therefore, $g(Z(\\breve{V})) \\subseteq Z(\\breve{V})$.\n\nHence, we have $g(Ut^{2}) \\subseteq g(Z(\\breve{V})) \\subseteq Z(\\breve{V}) = Vt^{2}$.\n\nNow we define a map $f : Vt + Ut^{2} + \\mu(V, V)t^{2} \\rightarrow Vt^{2}$ by\n$$f(x) =\n\\left\\{\n\\begin{aligned}\ng(x) \\cap Vt^{2},\\quad x \\in Vt\\\\\ng(x),\\quad x \\in Ut^{2}\\\\\n0,\\quad x \\in \\mu(V, V)t^{2}\n\\end{aligned}\n\\right.\n$$\nIt's clear that $f$ is linear.\n\nNote that\n\\[f(\\breve{\\mu}(\\breve{V}, \\breve{V})) = f(\\mu(V, V)t^{2}) = 0,\\]\n\\[\\breve{\\mu}(f(\\breve{V}), \\breve{\\alpha}^{k}(\\breve{V})) \\subseteq \\breve{\\mu}(Vt^{2}, \\alpha^{k}(V)t + \\alpha^{k}(V)t^{2}) = 0,\\]\nwe have $f \\in ZDer_{\\breve{\\alpha}^{k}}(\\breve{V})$.\n\nSince\n\\[(g - f)(Vt) = g(Vt) - f(Vt) = g(Vt) - g(Vt) \\cap Vt^{2} = g(Vt) - Vt^{2} \\subseteq Vt,\\]\n\\[(g - f)(Ut^{2}) = 0,\\]\n\\[(g - f)(\\mu(V, V)t^{2}) = g(\\breve{\\mu}(\\breve{V}, \\breve{V})) \\subseteq \\breve{\\mu}(\\breve{V}, \\breve{V}) = \\mu(V, V)t^{2},\\]\nhence there exist $D, D^{'} \\in End (V)$ such that $\\forall a \\in V, b \\in \\mu(V, V)$,\n\\[(g - f)(at) = D(a)t, (g - f)(bt^{2}) = D^{'}(b)t^{2}.\\]\n\nSince $g - f \\in Der_{\\breve{\\alpha}^{k}}(\\breve{V})$, we have\n\\[(g - f)(\\breve{\\mu}(a_{1}t, a_{2}t)) = \\breve{\\mu}((g - f)(a_{1}t), \\breve{\\alpha}^{k}(a_{2}t)) + \\breve{\\mu}(\\breve{\\alpha}^{k}(a_{1}t), (g - f)(a_{2}t)),\\]\nwhich implies that\n\\[D^{'}(\\mu(a_{1}, a_{2}))t^{2} = \\mu(D(a_{1}), \\alpha^{k}(a_{2}))t^{2} + \\mu(\\alpha^{k}(a_{1}), D(a_{2}))t^{2}.\\]\nHence, we have $D \\in QDer_{\\alpha^{k}}(V) \\subseteq QDer(V)$.\n\nTherefore, $g - f = \\varphi(D) \\in \\varphi(QDer(V))$, which implies that $Der(\\breve{V}) \\subseteq \\varphi(QDer(V)) + ZDer(\\breve{V})$. According to Proposition \\ref{prop:4.2} (2), we have $Der(\\breve{V}) = \\varphi(QDer(V)) + ZDer(\\breve{V})$.\n\n$\\forall f \\in \\varphi(QDer(V)) \\cap ZDer(\\breve{V})$, there exists $D \\in QDer(V)$ such that $f = \\varphi(D)$. For all $a \\in V, b \\in \\mu(V, V), u \\in U$,\n\\[f(at + bt^{2} + ut^{2}) = \\varphi(D)(at + bt^{2} + ut^{2}) = D(a)t + D^{'}(b)t^{2}.\\]\n\nOn the other hand, for any $xt + yt^{2} \\in \\breve{V}$, there exists $x^{'}t + y^{'}t^{2} \\in \\breve{V}$ such that $xt + yt^{2} = \\breve{\\alpha}^{k}(x^{'}t + y^{'}t^{2})$ since $\\breve{\\alpha}$ is surjective. Then we have\n\\begin{align*}\n&\\breve{\\mu}(f(at + bt^{2} + ut^{2}), xt + yt^{2}) = \\breve{\\mu}(f(at + bt^{2} + ut^{2}), \\breve{\\alpha}^{k}(x^{'}t + y^{'}t^{2}))\\\\\n&= f(\\breve{\\mu}(at + bt^{2} + ut^{2}, x^{'}t + y^{'}t^{2}))\\\\\n&= 0\n\\end{align*}\nsince $f \\in ZDer(\\breve{V})$. Hence, $f(at + bt^{2} + ut^{2}) \\in Z(\\breve{V}) = Vt^{2}$.\n\nTherefore, $D(a) = 0$, i.e., $D = 0$. Hence, $f = 0$.\n\nTherefore, $Der(\\breve{V}) = \\varphi(QDer(V)) \\dotplus ZDer(\\breve{V})$.\n\\end{proof}\n\n\\section{The centroids of Hom-Jordan algebras}\\label{se:5}\n\n\\begin{prop}\nSuppose that $(V, \\mu, \\alpha)$ is a simple multiplicative Hom-Jordan algebra over an algebraically closed field $\\rm{\\mathbb{F}}$ of characteristic $0$. If $C(V) = \\rm{\\mathbb{F}}id$, then $\\alpha = id$.\n\\end{prop}\n\\begin{proof}\nFor all $k \\in \\mathbb{N}^{+}$, $\\forall\\; 0 \\neq \\psi \\in C_{\\alpha^{k}}(V)$, we have $\\psi = p\\;id,\\quad p \\in \\rm{\\mathbb{F}},\\;p \\neq 0$. So $\\forall x, y \\in V$,\n\\[p\\mu(x, y) = \\psi(\\mu(x, y)) = \\mu(\\psi(x), \\alpha^{k}(y)) = \\mu(px, \\alpha^{k}(y)) = p\\mu(x, \\alpha^{k}(y)),\\]\nwhich implies that $\\mu(x, y) = \\mu(x, \\alpha^{k}(y))$ for all $k \\in \\mathbb{N}^{+}$.\n\nSince $\\rm{\\mathbb{F}}$ is algebraically closed, $\\alpha$ has an eigenvalue $\\lambda$. We denote the corresponding eigenspace by $E_{\\lambda}(\\alpha)$. So $E_{\\lambda}(\\alpha) \\neq 0$. Let $k = 1$. For any $x \\in E_{\\lambda}(\\alpha)$, $y \\in V$, we have\n\\[\\alpha(\\mu(x, y)) = \\mu(\\alpha(x), \\alpha(y)) = \\mu(\\lambda x, \\alpha(y)) = \\lambda\\mu(x, \\alpha(y)) = \\lambda\\mu(x, y),\\]\nwhich implies that $\\mu(x, y) \\in E_{\\lambda}(\\alpha)$.\n\nMoreover, for any $x \\in E_{\\lambda}(\\alpha)$, we have\n\\[\\alpha(\\alpha(x)) = \\lambda^{2}x = \\lambda\\alpha(x),\\]\nwhich implies that $\\alpha(x) \\in E_{\\lambda}(\\alpha)$, i.e., $\\alpha(E_{\\lambda}(\\alpha)) \\subseteq E_{\\lambda}(\\alpha)$.\n\nSo $E_{\\lambda}(\\alpha)$ is a Hom-ideal of $(V, \\mu, \\alpha)$. Since $(V, \\mu, \\alpha)$ is simple, we have $E_{\\lambda}(\\alpha) = V$, i.e., $\\alpha = \\lambda id$.\n\nThen for any $x, y \\in V, k = 1$, we have\n\\[\\mu(x, y) = \\mu(x, \\alpha(y)) = \\mu(x, \\lambda y) = \\lambda\\mu(x, y),\\]\nwhich implies that $\\lambda = 1$, i.e., $\\alpha = id$.\n\\end{proof}\n\\begin{prop}\nLet $(V, \\mu, \\alpha)$ be a multiplicative Hom-Jordan algebra.\n\\begin{enumerate}[(1)]\n\\item If $\\alpha$ is a surjection, then $V$ is indecomposable if and only if $C(V)$ does not contain idempotents except $0$ and $id$.\n\\item If $(V, \\mu, \\alpha)$ is perfect(i.e., $V = \\mu(V, V)$), then every $\\psi \\in C_{\\alpha^{k}}(V)$ is $\\alpha^{k}$-symmetric with respect to any invariant form on $V$.\n\\end{enumerate}\n\\end{prop}\n\\begin{proof}\n(1). $(\\Rightarrow)$. If there exists $\\psi \\in C_{\\alpha^{k}}(V)$ is an idempotent and satisfies that $\\psi \\neq 0,\\;id$, then $\\psi^{2}(x) = \\psi(x),\\quad x \\in V$.\n\nFor any $x \\in ker(\\psi), y \\in V$, we have\n\\[\\psi(\\mu(x, y)) = \\mu(\\psi(x), \\alpha^{k}(y)) = \\mu(0, \\alpha^{k}(y)) = 0,\\]\nwhich implies that $\\mu(x, y) \\in ker(\\psi)$.\nMoreover, $\\psi(\\alpha(x)) = \\alpha(\\psi(x)) = \\alpha(0) = 0$, so $\\alpha(x) \\in ker(\\psi)$, i.e., $\\alpha(ker(\\psi)) \\subseteq ker(\\psi)$.\n\nHence, $ker(\\psi)$ is a Hom-ideal of $(V, \\mu, \\alpha)$.\n\nFor any $y \\in Im(\\psi)$, there exists $y^{'} \\in V$ such that $y = \\psi(y^{'})$. And for any $z \\in V$, there exists $z^{'} \\in V$ such that $z = \\alpha^{k}(z^{'})$ since $\\alpha$ is surjective. Then\n\\[\\mu(y, z) = \\mu(\\psi(y^{'}), \\alpha^{k}(z^{'})) = \\psi(\\mu(y^{'}, z^{'})),\\]\nwhich implies that $\\mu(y, z) \\in Im(\\psi)$.\nMoreover, $\\alpha(y) = \\alpha(\\psi(y^{'})) = \\psi(\\alpha(y^{'}))$, so $\\alpha(y) \\in Im(\\psi)$, i.e., $\\alpha(Im(\\psi)) \\subseteq Im(\\psi)$.\n\nHence, $Im(\\psi)$ is a Hom-ideal of $(V, \\mu, \\alpha)$.\n\n$\\forall x \\in ker(\\psi) \\cap Im(\\psi)$, $\\exists x^{'} \\in V$ such that $x = \\psi(x^{'})$. So $0 = \\psi(x) = \\psi^{2}(y) = \\psi(y)$, which implies that $x = 0$. Hence, $ker(\\psi) \\cap Im(\\psi) = \\{0\\}$.\n\nWe have a decomposition $x = (x - \\psi(x)) + \\psi(x)$, $\\forall x \\in V$. So we have $V = ker(\\psi) \\oplus Im(\\psi)$. Contradiction.\n\n$(\\Leftarrow)$. Suppose that $V = V_{1} \\oplus V_{2}$ where $V_{i}$ are Hom-ideals of $(V, \\mu, \\alpha)$. Then for any $x \\in V$, $\\exists x_{i} \\in V_{i}$ such that $x = x_{1} + x_{2}$.\n\nDefine $\\psi : V \\rightarrow V$ by $\\psi(x) = \\psi(x_{1} + x_{2}) = x_{1} - x_{2}$. It's obvious that $\\psi^{2}(x) = \\psi(x)$.\n\nFor any $x, y \\in V$, suppose that $x = x_{1} + x_{2},\\;y = y_{1} + y_{2}$,\n\\[\\psi(\\mu(x, y)) = \\psi(\\mu(x_{1} + x_{2}, y_{1} + y_{2})) = \\psi(\\mu(x_{1}, y_{1}) + \\mu(x_{2}, y_{2})) = \\mu(x_{1}, y_{1}) - \\mu(x_{2}, y_{2}),\\]\n\\[\\mu(\\psi(x), y) = \\mu(\\psi(x_{1} + x_{2}), y_{1} + y_{2}) = \\mu(x_{1} - x_{2}, y_{1} + y_{2}) = \\mu(x_{1}, y_{1}) - \\mu(x_{2}, y_{2}),\\]\nwe have\n\\[\\psi(\\mu(x, y)) = \\mu(\\psi(x), y).\\]\nSimilarly, we have $\\psi(\\mu(x, y)) = \\mu(x, \\psi(y))$. So $\\psi \\in C_{\\alpha^{0}}(V) \\subseteq C(V)$. Contradiction.\n\n(2). Let $f$ be an invariant $\\rm{F}$-bilinear form on $V$. Then we have $f(\\mu(x, y), z) = f(x, \\mu(y, z))$ for any $x, y, z \\in V$.\n\nSince $(V, \\mu, \\alpha)$ is perfect,let $\\psi \\in C_{\\alpha^{k}}(V)$, then for any $a, b, c \\in V$ we have\n\\begin{align*}\n&f(\\psi(\\mu(a, b)), \\alpha^{k}(c)) = f(\\mu(\\alpha^{k}(a), \\psi(b)), \\alpha^{k}(c)) = f(\\alpha^{k}(a), \\mu(\\psi(b), \\alpha^{k}(c)))\\\\\n&= f(\\alpha^{k}(a), \\mu(\\alpha^{k}(b), \\psi(c))) = f(\\mu(\\alpha^{k}(a), \\alpha^{k}(b)), \\psi(c)) = f(\\alpha^{k}(\\mu(a, b)), \\psi(c)).\n\\end{align*}\n\\end{proof}\n\\begin{prop}\nLet $(V, \\mu, \\alpha)$ be a Hom-Jordan algebra and $I$ an $\\alpha$-invariant subspace of $V$ where $\\alpha|_{I}$ is surjective. Then $Z_{V}(I) = \\{x \\in V | \\mu(x, y) = 0,\\;\\forall y \\in I\\}$ is invariant under $C(V)$. So is any perfect Hom-ideal of $(V, \\mu, \\alpha)$.\n\\end{prop}\n\\begin{proof}\nFor any $\\psi \\in C_{\\alpha^{k}}(V)$ and $x \\in Z_{V}(I)$, $\\forall y \\in I$, $\\exists y^{'} \\in I$ such that $y = \\alpha^{k}(y^{'})$ since $\\alpha|_{I}$ is surjective, we have\n\\[\\mu(\\psi(x), y) = \\mu(\\psi(x), \\alpha^{k}(y^{'})) = \\psi(\\mu(x, y^{'})) = \\psi(0) = 0,\\]\nwhich implies that $\\psi(x) \\in Z_{V}(I)$. So $Z_{V}(I)$ is invariant under $C(V)$.\n\nSuppose that $J$ is a perfect Hom-ideal of $(V, \\mu, \\alpha)$. Then $J = \\mu(J, J)$. For any $y \\in J$, there exist $a, b \\in J$ such that $y = \\mu(a, b)$, then we have\n\\[\\psi(y) = \\psi(\\mu(a, b)) = \\mu(\\psi(a), \\alpha^{k}(b)) \\in \\mu(J, J) = J.\\]\nSo $J$ is invariant under $C(V)$.\n\\end{proof}\n\\begin{thm}\nSuppose that $(V_{1}, \\mu_{1}, \\alpha_{1})$ and $(V_{2}, \\mu_{2}, \\alpha_{2})$ are two Hom-Jordan algebras over field $\\rm{F}$ with $\\alpha_{1}$ is surjective. Let $\\pi : V_{1} \\rightarrow V_{2}$ be an epimorphism of Hom-Jordan algebras. Then for any $f \\in End_{\\rm{F}}(V_{1}, ker(\\pi)) := \\{g \\in \\mathcal{W}_{1} | g(ker(\\pi)) \\subseteq ker(\\pi)\\}$, there exists a unique $\\bar{f} \\in \\mathcal{W}_{2}$ satisfying $\\pi \\circ f = \\bar{f} \\circ \\pi$ where $\\mathcal{W}_{i}(i = 1, 2)$ are defined as Lemma \\ref{le:2.7}. Moreover, the following results hold:\n\\begin{enumerate}[(1)]\n\\item The map $\\pi_{End} : End_{\\rm{F}}(V_{1}, ker(\\pi)) \\rightarrow \\mathcal{W}_{2}$, $f \\mapsto \\bar{f}$ is a Hom-algebra homomorphism with the following proporties, where $(End_{\\rm{F}}(V_{1}, ker(\\pi)), \\circ, \\sigma_{1})$ and $(\\mathcal{W}_{2}, \\circ, \\sigma_{2})$ are two Hom-associative algebras and $\\sigma_{1} : End_{\\rm{F}}(V_{1}, ker(\\pi)) \\rightarrow End_{\\rm{F}}(V_{1}, ker(\\pi)),\\;f \\mapsto \\alpha_{1} \\circ f$ and $\\sigma_{2} : \\mathcal{W}_{2} \\rightarrow \\mathcal{W}_{2},\\;g \\mapsto \\alpha_{2} \\circ g$.\n \\begin{enumerate}[(a)]\n \\item $\\pi_{End}(Mult(V_{1})) = Mult(V_{2})$, $\\pi_{End}(C(V_{1}) \\cap End_{\\rm{F}}(V_{1}, ker(\\pi))) \\subseteq C(V_{2})$.\n \\item By restriction, there is a Hom-algebra homomorphism\n \\[\\pi_{C} : C(V_{1}) \\cap End_{\\rm{F}}(V_{1}, ker(\\pi)) \\rightarrow C(V_{2}),\\quad f \\mapsto \\bar{f}.\\]\n \\item If $ker(\\pi) = Z(V_{1})$, then every $\\varphi \\in C(V_{1})$ leaves $ker(\\pi)$ invariant.\n \\end{enumerate}\n\\item Suppose that $V_{1}$ is perfect and $ker(\\pi) \\subseteq Z(V_{1})$. Then $\\pi_{C} : C(V_{1}) \\cap End_{\\rm{F}}(V_{1}, ker(\\pi)) \\rightarrow C(V_{2}),\\quad f \\mapsto \\bar{f}$ is injective.\n\\item If $V_{1}$ is perfect, $Z(V_{2}) = \\{0\\}$ and $ker(\\pi) \\subseteq Z(V_{1})$, then $\\pi_{C} : C(V_{1}) \\rightarrow C(V_{2})$ is a Hom-algebra homomorphism.\n\\end{enumerate}\n\\end{thm}\n\\begin{proof}\nFor any $y \\in V_{2}$, there exists $x \\in V_{1}$ such that $y = \\pi(x)$ since $\\pi$ is an epimorphism. Then for any $f \\in End_{\\rm{F}}(V_{1}, ker(\\pi))$, define $\\bar{f} : V_{2} \\rightarrow V_{2}$ by $\\bar{f}(y) = \\pi(f(x))$ where $x$ satisfies that $y = \\pi(x)$. It's obvious that $\\pi \\circ f = \\bar{f} \\circ \\pi$.\n\\begin{align*}\n&(\\alpha_{2} \\circ \\bar{f}) \\circ \\pi = \\alpha_{2} \\circ (\\bar{f} \\circ \\pi) = \\alpha_{2} \\circ (\\pi \\circ f) = (\\alpha_{2} \\circ \\pi) \\circ f = (\\pi \\circ \\alpha_{1}) \\circ f = \\pi \\circ (\\alpha_{1} \\circ f)\\\\\n&= \\pi \\circ (f \\circ \\alpha_{1}) = (\\pi \\circ f) \\circ \\alpha_{1} = (\\bar{f} \\circ \\pi) \\circ \\alpha_{1} = \\bar{f} \\circ (\\pi \\circ \\alpha_{1}) = \\bar{f} \\circ (\\alpha_{2} \\circ \\pi) = (\\bar{f} \\circ \\alpha_{2}) \\circ \\pi.\n\\end{align*}\nSince $\\pi$ is surjective, we have $\\alpha_{2} \\circ \\bar{f} = \\bar{f} \\circ \\alpha_{2}$, i.e., $\\bar{f} \\in \\mathcal{W}_{2}$.\n\nIf there exist $\\bar{f}$ and $\\bar{f}^{'}$ such that $\\pi \\circ f = \\bar{f} \\circ \\pi$ and $\\pi \\circ f = \\bar{f}^{'} \\circ \\pi$, then we have $\\bar{f} \\circ \\pi = \\bar{f}^{'} \\circ \\pi$. For any $y \\in V_{2}$, there exists $x \\in V_{1}$ such that $y = \\pi(x)$. Hence, $\\bar{f}(y) = \\bar{f}(\\pi(x)) = \\bar{f}^{'}(\\pi(x)) = \\bar{f}^{'}(y)$. So $\\bar{f} = \\bar{f}^{'}$.\n\n(1). For all $f, g \\in End_{\\rm{F}}(V_{1}, ker(\\pi))$, we have\n\\[\\pi \\circ (f \\circ g) = (\\pi \\circ f) \\circ g = (\\bar{f} \\circ \\pi) \\circ g = \\bar{f} \\circ (\\pi \\circ g) = \\bar{f} \\circ (\\bar{g} \\circ \\pi) = (\\bar{f} \\circ \\bar{g}) \\circ \\pi,\\]\nwhich implies that $\\pi_{End}(f \\circ g) = \\bar{f} \\circ \\bar{g} = \\pi_{End}(f) \\circ \\pi_{End}(g)$.\n\\[\\pi \\circ (\\alpha_{1} \\circ f) = (\\pi \\circ \\alpha_{1}) \\circ f = (\\alpha_{2} \\circ \\pi) \\circ f = \\alpha_{2} \\circ (\\pi \\circ f) = \\alpha_{2} \\circ (\\bar{f} \\circ \\pi) = (\\alpha_{2} \\circ \\bar{f}) \\circ \\pi,\\]\nwhich implies that $\\pi_{End}(\\sigma_{1}(f)) = \\alpha_{2} \\circ \\bar{f} = \\sigma_{2}(\\pi_{End}(f))$, i.e., $\\pi_{End} \\circ \\sigma_{1} = \\sigma_{2} \\circ \\pi_{End}$. Therefore, $\\pi_{End}$ is a Hom-algebra homomorphism.\n\n(a). For all $x \\in V_{1}$, $z \\in ker(\\pi)$, we have\n\\[\\pi(L_{x}(z)) = \\pi(\\mu_{1}(x, z)) = \\mu_{2}(\\pi(x), \\pi(z)) = \\mu_{2}(\\pi(x), 0) = 0,\\]\nwhich implies that $L_{x}(z) \\in ker(\\pi)$, i.e., $L_{x} \\in End_{\\rm{F}}(V_{1}, ker(\\pi))$.\n\nMoreover, we have $\\pi \\circ L_{x} = L_{\\pi(x)} \\circ \\pi$. So $\\pi_{End}(L_{x}) = L_{\\pi(x)} \\in Mult(V_{2})$. Hence, $\\pi_{End}(Mult(V_{1})) \\subseteq Mult(V_{2})$.\n\nOn the other hand, for any $L_{y} \\in Mult(V_{2})$, $\\exists x \\in V_{1}$ such that $y = \\pi(x)$. So\n\\[L_{y} = L_{\\pi(x)} = \\pi_{End}(L_{x}) \\in \\pi_{End}(Mult(V_{1})),\\]\nwhich implies that $Mult(V_{2}) \\subseteq \\pi_{End}(Mult(V_{1}))$. Therefore, $\\pi_{End}(Mult(V_{1})) = Mult(V_{2})$.\n\nFor any $\\varphi \\in C_{\\alpha_{1}^{k}}(V_{1}) \\cap End_{\\rm{F}}(V_{1}, ker(\\pi))$, $\\forall x^{'}, y^{'}\\in V_{2}$, $\\exists x, y \\in V_{1}$ such that $x^{'} = \\pi(x),\\; y^{'} = \\pi(y)$, then we have\n\\begin{align*}\n&\\bar{\\varphi}(\\mu_{2}(x^{'}, y^{'})) = \\bar{\\varphi}(\\mu_{2}(\\pi(x), \\pi(y))) = \\bar{\\varphi}(\\pi(\\mu_{1}(x, y))) = \\pi(\\varphi(\\mu_{1}(x, y))) = \\pi(\\mu_{1}(\\varphi(x), \\alpha_{1}^{k}(y)))\\\\\n&= \\mu_{2}(\\pi(\\varphi(x)), \\pi(\\alpha_{1}^{k}(y))) = \\mu_{2}(\\bar{\\varphi}(\\pi(x)), \\alpha_{2}^{k}(\\pi(y))) = \\mu_{2}(\\bar{\\varphi}(x^{'}), \\alpha_{2}^{k}(y^{'})).\n\\end{align*}\nSimilarly, we have $\\bar{\\varphi}(\\mu_{2}(x^{'}, y^{'})) = \\mu_{2}(\\alpha_{2}^{k}(x^{'}), \\bar{\\varphi}(y^{'}))$. Hence, $\\bar{\\varphi} \\in C_{\\alpha_{2}^{k}}(V_{2})$. Therefore, $\\pi_{End}(C(V_{1}) \\cap End_{\\rm{F}}(V_{1}, ker(\\pi))) \\subseteq C(V_{2})$.\n\n(b). $\\forall f, g \\in C(V_{1}) \\cap End_{\\rm{F}}(V_{1}, ker(\\pi))$, we have\n\\[\\pi \\circ (f \\circ g) = (\\pi \\circ f) \\circ g = (\\bar{f} \\circ \\pi) \\circ g = \\bar{f} \\circ (\\pi \\circ g) = \\bar{f} \\circ (\\bar{g} \\circ \\pi) = (\\bar{f} \\circ \\bar{g}) \\circ \\pi,\\]\nwhich implies that $\\pi_{C}(f \\circ g) = \\bar{f} \\circ \\bar{g} = \\pi_{C}(f) \\circ \\pi_{C}(g)$.\n\\[\\pi \\circ (\\alpha_{1} \\circ f) = (\\pi \\circ \\alpha_{1}) \\circ f = (\\alpha_{2} \\circ \\pi) \\circ f = \\alpha_{2} \\circ (\\pi \\circ f) = \\alpha_{2} \\circ (\\bar{f} \\circ \\pi) = (\\alpha_{2} \\circ \\bar{f}) \\circ \\pi,\\]\nwhich implies that $\\pi_{C}(\\sigma_{1}(f)) = \\alpha_{2} \\circ \\bar{f} = \\sigma_{2}(\\pi_{C}(f))$, i.e., $\\pi_{C} \\circ \\sigma_{1} = \\sigma_{2} \\circ \\pi_{C}$. Therefore, $\\pi_{C}$ is a Hom-algebra homomorphism.\n\n(c). If $ker(\\pi) = Z(V_{1})$, $\\forall x \\in ker(\\pi), \\varphi \\in C_{\\alpha_{1}^{k}}(V_{1})$, for any $y \\in V_{1}$ there exists $y^{'} \\in V_{1}$ such that $y = \\alpha_{1}^{k}(y^{'})$ since $\\alpha_{1}$ is surjective. We have\n\\[\\mu_{1}(\\varphi(x), y) = \\mu_{1}(\\varphi(x), \\alpha_{1}^{k}(y^{'})) = \\varphi(\\mu_{1}(x, y^{'})) = \\varphi(0) = 0,\\]\nwhich implies that $\\varphi(x) \\in Z(V_{1}) = ker(\\pi)$. Hence, every $\\varphi \\in C(V_{1})$ leaves $ker(\\pi)$ invariant.\n\n(2). If $\\bar{\\varphi} = 0$ for $\\varphi \\in C_{\\alpha_{1}^{k}}(V_{1}) \\cap End_{\\rm{F}}(V_{1}, ker(\\pi))$, then $\\pi(\\varphi(V_{1})) = \\bar{\\varphi}(\\pi(V_{1})) = 0$. Hence, $\\varphi(V_{1}) \\subseteq ker(\\pi) \\subseteq Z(V_{1})$.\n\nHence, $\\varphi(\\mu_{1}(x, y)) = \\mu_{1}(\\varphi(x), \\alpha_{1}^{k}(y)) = 0$. Note that $V_{1}$ is perfect, we have $\\varphi = 0$. Therefore, $\\pi_{C} : C(V_{1}) \\cap End_{\\rm{F}}(V_{1}, ker(\\pi)) \\rightarrow C(V_{2}),\\quad f \\mapsto \\bar{f}$ is injective.\n\n(3). $\\forall y \\in V_{2}, \\exists y^{'} \\in V_{1}$ such that $y = \\pi(y^{'})$. For all $x \\in Z(V_{1})$,\n\\[\\mu_{2}(\\pi(x), y) = \\mu_{2}(\\pi(x), \\pi(y^{'})) = \\pi(\\mu_{1}(x, y^{'})) = \\pi(0) = 0,\\]\nwhich implies that $\\pi(x) \\in Z(V_{2})$. So $\\pi(Z(V_{1})) \\subseteq Z(V_{2}) = \\{0\\}$. Therefore, $Z(V_{1}) \\subseteq ker(\\pi)$. Note that $ker(\\pi) \\subseteq Z(V_{1})$, we have $Z(V_{1}) = ker(\\pi)$.\n\nFor all $\\varphi \\in C_{\\alpha_{1}^{k}}(V_{1}), x \\in Z(V_{1}), y \\in V_{1}$, $\\exists y^{'} \\in V_{1}$ such that $y = \\alpha_{1}^{k}(y^{'})$,\n\\[\\mu_{1}(\\varphi(x), y) = \\mu_{1}(\\varphi(x), \\alpha_{1}^{k}(y^{'})) = \\varphi(\\mu_{1}(x, y^{'})) = \\varphi(0) = 0,\\]\nwhich implies that $\\varphi(x) \\in Z(V_{1}) = ker(\\pi)$. So $\\varphi(ker(\\pi)) \\subseteq ker(\\pi)$, i.e., $\\varphi \\in End_{\\rm{F}}(V_{1}, ker(\\pi))$. So $C(V_{1}) \\subseteq End_{\\rm{F}}(V_{1}, ker(\\pi))$. Therefore, $C(V_{1}) \\cap End_{\\rm{F}}(V_{1}, ker(\\pi)) = C(V_{1})$. According to (1) (b), we have $\\pi_{C} : C(V_{1}) \\rightarrow C(V_{2})$ is a Hom-algebra homomorphism.\n\\end{proof}\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\n\\section{Introduction}\n\\label{intro}\n\nMeasurements of single-top-quark production via the electroweak interaction, a process first observed in proton--antiproton (\\ensuremath{{p}\\bar{p}}\\xspace) collisions at the Tevatron~\\cite{Aaltonen:2009jj, Abazov:2009ii}, have entered the precision era at the Large Hadron Collider (LHC). It has become possible to measure top-quark properties using single-top-quark events~\\cite{Giammanco:2017xyn}. Single-top-quark production is sensitive to new physics mechanisms~\\cite{Tait:2000sh} that either modify the \\ensuremath{tWb}\\xspace coupling~\\cite{Cao:2007ea,Godbole:2011vw,Zhang:2010dr,AguilarSaavedra:2008gt,Dror:2015nkp,Berger:2015yha} or introduce new particles and interactions~\\cite{He:1999vp,Aguilar-Saavedra:2013qpa,Durieux:2014xla,Nutter:2012an,Hashemi:2013raa,Drueke:2014pla}. The production rate of single top quarks is proportional to the square of the left-handed coupling at the \\ensuremath{tWb}\\xspace production vertex, assuming that there are no significant \\ensuremath{tWd}\\xspace or \\ensuremath{tWs}\\xspace contributions. In the Standard Model (SM), this coupling is given by the Cabibbo--Kobayashi--Maskawa (CKM)~\\cite{CKM1,CKM2} matrix element $\\ensuremath{V_{\\myt\\myb}}\\xspace$. Indirect measurements of $|\\ensuremath{V_{\\myt\\myb}}\\xspace|$, from precision measurements of $B$-meson decays~\\cite{Tanabashi:2018oca} and from top-quark decays~\\cite{Abazov:2011zk,Aaltonen:2013luz,Aaltonen:2014yua,CMS-TOP-12-035}, rely on the SM assumptions that the CKM matrix is unitary\nand that there are three quark generations. The most stringent indirect determination comes from a global fit to all available $B$-physics measurements, resulting in $|\\ensuremath{V_{\\myt\\myb}}\\xspace| = 0.999105 \\pm 0.000032$~\\cite{Tanabashi:2018oca}. This fit also assumes the absence of any new physics mechanisms that might affect \\bquarks. The most precise indirect measurement using top-quark events was performed by the CMS Collaboration in proton--proton (\\ensuremath{pp}\\xspace) collisions at a centre-of-mass energy of \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, resulting in $|\\ensuremath{V_{\\myt\\myb}}\\xspace| = 1.007\\pm0.016$~\\cite{CMS-TOP-12-035}.\n\nA direct estimate of the coupling at the \\ensuremath{tWb}\\xspace production vertex, $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$, is obtained from the measured single-top-quark cross-section $\\ensuremath{\\sigma_{\\rm meas.}}\\xspace$ and its corresponding theoretical expectation $\\ensuremath{\\sigma_{\\rm theo.}}\\xspace$,\n\\begin{eqnarray}\n\\label{eq:vtb}\n\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace = \\sqrt{\\frac{\\ensuremath{\\sigma_{\\rm meas.}}\\xspace}{\\ensuremath{\\sigma_{\\rm theo.}}\\xspace \\mathsmaller{\\;(\\ensuremath{V_{\\myt\\myb}}\\xspace=1)}}}.\n\\end{eqnarray}\nThe $\\ensuremath{f_{\\rm LV}}\\xspace$ term is a form factor, assumed to be real, that parameterises the possible presence of anomalous left-handed vector couplings~\\cite{AguilarSaavedra:2008zc}. By construction, this form factor is exactly one in the SM, while it can be different from one in models of new physics processes. The direct estimation assumes that $|\\ensuremath{V_{td}}\\xspace|,|\\ensuremath{V_{ts}}\\xspace| \\ll |\\ensuremath{V_{\\myt\\myb}}\\xspace|$~\\cite{Abazov:2012uga,Alwall:2006bx}, and that the \\ensuremath{tWb}\\xspace interaction involves a left-handed weak coupling, like that in the SM. The $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ determination via single-top-quark production is independent of assumptions about the number of quark generations and the unitarity of the CKM matrix~\\cite{Tait:1999cf,Belyaev:2000me,Tait:2000sh,Cao:2015qta}. Since the indirect determination of $|\\ensuremath{V_{\\myt\\myb}}\\xspace|$ gives a value close to unity, $\\ensuremath{V_{\\myt\\myb}}\\xspace$ is considered equal to one in theoretical calculations of the single-top-quark cross-section. The combination of single-top-quark measurements from the Tevatron gives $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace=1.02_{-0.05}^{+0.06}$~\\cite{Aaltonen:2015cra}.\n\nSingle-top-quark production at a hadron collider mostly proceeds, according to the SM prediction, via three modes that can be defined at leading order (LO) in perturbative quantum chromodynamics (QCD): the exchange of a virtual \\ensuremath{W}\\xspace boson in the \\myt-channel\\xspace\\ or in the \\mys-channel\\xspace, and the associated production of a top quark and a \\ensuremath{W}\\xspace boson (\\ensuremath{tW}\\xspace). Representative Feynman diagrams for these processes at LO are shown in Figure~\\ref{fig:feyn}.\n\n\\begin{figure}[!h!tbp]\n \\begin{center}\n \\subfloat[]{ \\includegraphics[height=0.12\\textheight]{intro\/tchanFeyn} }\n \\subfloat[]{ \\includegraphics[height=0.12\\textheight]{intro\/WtFeyn} }\n \\subfloat[]{ \\includegraphics[height=0.12\\textheight]{intro\/schanFeyn} }\n \\caption{Representative Feynman diagrams at LO in QCD and in the five-flavour-number scheme for single-top-quark production in (a) the \\myt-channel\\xspace, (b) \\ensuremath{tW}\\xspace production, and (c) the \\mys-channel\\xspace.}\n \\label{fig:feyn}\n \\end{center}\n\\end{figure}\n\nIn \\ensuremath{pp}\\xspace collisions at the LHC, the process with the largest single-top-quark production cross-section is the \\myt-channel\\xspace, where a light-flavour quark \\ensuremath{q}\\xspace from one of the colliding protons interacts with a \\bquark by exchanging a space-like virtual \\ensuremath{W}\\xspace boson, producing a top quark (\\ensuremath{t}\\xspace-quark) and a recoiling light-flavour quark $\\ensuremath{q}\\xspace'$, called the spectator quark.\nFor \\myt-channel\\xspace production at LO, the \\bquark can be considered as directly emitted from the other proton (five-flavour-number scheme or 5FS) or it can come from gluon splitting (four-flavour-number scheme or 4FS)~\\cite{Frederix:2012dh}. The kinematic properties of the spectator quark provide distinctive features for this process~\\cite{TOPQ-2015-05,CMS-TOP-12-038}. The associated production of a \\ensuremath{W}\\xspace boson and a top quark has the second-largest production cross-section. In a representative process of \\ensuremath{tW}\\xspace production, a gluon interacts with an initial \\bquark by exchanging a virtual \\bquark, producing a \\ensuremath{t}\\xspace-quark and a \\ensuremath{W}\\xspace boson. The measurement of this process suffers from a large background from top-quark pair (\\ensuremath{t \\bar{t}}\\xspace) production~\\cite{TOPQ-2012-20,CMS-TOP-12-040}. The \\mys-channel\\xspace\\ cross-section is the smallest at the LHC. In this process, a quark--antiquark pair annihilates to produce a time-like virtual \\ensuremath{W}\\xspace boson, which decays to a \\ensuremath{t}\\xspace-quark and a \\ensuremath{\\bar{b}}\\xspace-quark. This process was observed in \\ensuremath{{p}\\bar{p}}\\xspace collisions at the Tevatron~\\cite{CDF:2014uma} \nand evidence of it was reported by the ATLAS Collaboration in \\ensuremath{pp}\\xspace collisions at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}~\\cite{TOPQ-2015-01}.\n\nIn this paper, the \\myt-channel\\xspace, \\ensuremath{tW}\\xspace, and \\mys-channel\\xspace\\ single-top-quark cross-section measurements by the ATLAS and CMS experiments are combined for each production mode, separately at $pp$ centre-of-mass energies of \\ensuremath{ 7 \\text{ and } 8~\\TeV}. A combined determination of $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ is also presented, using as inputs the values of $\\ensuremath{\\absvtb^2}\\xspace$ calculated from the measured and predicted single-top-quark cross-sections in the three production modes at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}. Using the same approach, results are also shown for $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ combinations for each production mode.\n\nThe theoretical cross-section calculations are described in Section~\\ref{sec:theory}. Section~\\ref{sec:atlascmsxs} presents the cross-section measurements. The combination methodology is briefly described in Section~\\ref{sec:method}. Section~\\ref{sec:systcat} is devoted to a discussion of systematic uncertainties in the cross-section measurements as well as theoretical calculations, where the latter affect the $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ extraction in particular.\nThe assumptions made about the correlation of uncertainties between the two experiments, as well as between theoretical calculations, are also discussed. \nSection~\\ref{sec:xscomb} presents the combination of cross-sections for each production mode at the same centre-of-mass energy.\nIn Section~\\ref{sec:vtbcomb}, determinations of $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ are performed using all single-top-quark cross-section measurements together or by production mode. Stability tests are also shown and discussed. In Section~\\ref{sec:sum}, the results are summarised.\n\n\n\\section{Theoretical cross-section calculations}\n\\label{sec:theory}\n\nThe theoretical predictions for the single-top-quark production cross-sections are calculated at next-to-leading order (NLO) in the strong coupling constant $\\alpha_{\\rm s}$, at NLO with next-to-next-to-leading-logarithm (NNLL) resummation (named NLO+NNLL), and at next-to-next-to-leading order (NNLO). The difference between 4FS and 5FS is small~\\cite{Campbell:2009ss, Campbell:2009gj}, and the calculations use the 5FS.\nThe NLO prediction is used in the \\ensuremath{V_{\\myt\\myb}}\\xspace\\ combination for the \\myt-channel\\xspace\\ and \\mys-channel\\xspace, while the NLO+NNLL prediction is used for \\ensuremath{tW}\\xspace, as explained below.\nThe NLO prediction is calculated with \\textsc{HatHor}\\xspace (v2.1)~\\cite{Aliev:2010zk,Kant:2014oha}. Uncertainties comprise the scale uncertainty, the $\\alpha_{\\rm s}$ uncertainty, and the parton distribution function (PDF) uncertainty. The scale uncertainty is evaluated by varying the renormalisation and factorisation scales up and down together by a factor of two. The combination of the PDF+$\\alpha_{\\rm s}$ uncertainty is calculated according to the PDF4LHC prescription~\\cite{Botje:2011sn} from the envelope of the uncertainties at 68\\% confidence level (CL) in the MSTW2008 NLO, CT10 NLO~\\cite{Lai:2010vv}, and NNPDF2.3~\\cite{Ball:2012cx} PDF sets. \n\nThe NLO+NNLL predictions~\\cite{Kidonakis:2012rm} are available for all single-top-quark production modes~\\cite{Kidonakis:2011wy, Kidonakis:2010ux, Kidonakis:2010tc}. Uncertainties in these calculations are estimated by varying the renormalisation and factorisation scales between $m_{\\myt}\\xspace\/2$ and $2 m_{\\myt}\\xspace$, where $m_{\\myt}\\xspace$ is the top-quark mass, and from the 90\\% CL uncertainties in the MSTW2008 NNLO~\\cite{MSTW, MSTW2008} PDF set. The evaluation of the PDF uncertainties is provided by the author of Refs.~\\cite{Kidonakis:2011wy, Kidonakis:2010ux, Kidonakis:2010tc} and is not fully compatible with the PDF4LHC prescription.\nThe \\myt-channel\\xspace\\ cross-sections at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ are also computed at NNLO in $\\alpha_{\\rm s}$~\\cite{Brucherseifer:2014ama}, with the renormalisation and factorisation scales set to $m_{\\myt}\\xspace$. This results in cross-sections which are about 0.3\\% and 0.6\\% lower than the NLO values at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ respectively. However, only a limited number of scale variations are evaluated~\\cite{Brucherseifer:2014ama}.\n\nA summary of all the available theoretical cross-section predictions for \\myt-channel\\xspace, \\ensuremath{tW}\\xspace, and \\mys-channel\\xspace production, $\\ensuremath{\\sigma^{\\myt\\textrm{-chan.}}_{\\mathrm{\\textrm{theo.}}}}\\xspace$, $\\ensuremath{\\sigma^{\\mytW}_{\\mathrm{\\textrm{theo.}}}}\\xspace$, and $\\ensuremath{\\sigma^{\\mys\\textrm{-chan.}}_{\\mathrm{\\textrm{theo.}}}}\\xspace$ respectively, with their uncertainties is shown in Table~\\ref{tab:xs}.\n\n\\begin{table}[!htb]\n \\centering\n \\caption{Predicted cross-sections for single-top-quark production at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ at the LHC. Uncertainties include scale and PDF+$\\alpha_{\\rm s}$ variations, except for the NNLO predictions, which only contain the scale variation. The PDF+$\\alpha_{\\rm s}$ uncertainties are evaluated according to the PDF4LHC prescription only for the NLO predictions. The uncertainties associated with the top-quark mass $m_{\\myt}\\xspace$ and beam energy $E_{\\rm beam}$ are also given for the NLO predictions for the $t$- and $s$-channels, and for the NLO+NNLL prediction for \\ensuremath{tW}\\xspace\\ production. The value of $m_{\\myt}\\xspace$ is set to $172.5~\\GeV$ in all predictions. The cross-sections marked with $^\\dagger$ are those used in the \\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace\\ combination.}\n \\IfPackagesLoaded{adjustbox}{\\begin{adjustbox}{max width=1.0\\textwidth}}{}\n \\begin{tabular}{l|l|l|l}\n \\hline\n \\hline\n \\T\\B\n $\\sqrt{s}$ & Process & Accuracy & $\\ensuremath{\\sigma_{\\rm theo.}}\\xspace{}$\\,[pb] \\\\\n \\hline\n \\T\\B\n & & NLO${}^{\\dagger}$ & $63.9_{-1.3}^{+1.9}~\\text{(scale)} \\pm 2.2~\\text{(PDF+$\\alpha_{\\rm s}$)} \\pm 0.7~\\text{($m_\\ensuremath{t}\\xspace$)} \\pm 0.1~\\text{($E_{\\rm beam}$)}$ \\\\\n \\cline{3-4}\n \\T\\B\n & \\myt-channel\\xspace & NLO+NNLL & $64.6^{+2.6}_{-1.7}~\\text{(scale+PDF+$\\alpha_{\\rm s}$)}$\\\\\n \\cline{3-4}\n \\T\\B\n & & NNLO & $63.7^{+0.5}_{-0.3}~\\text{(scale)}$\\\\\n \\cline{2-4}\n \\T\\B\n \\ensuremath{ 7~\\TeV} & \\multirow{2}{*}{\\ensuremath{tW}\\xspace} & NLO & $13.2_{-0.6}^{+0.5}~\\text{(scale)} \\pm 1.3~\\text{(PDF+$\\alpha_{\\rm s}$)}$ \\\\\n \\cline{3-4}\n \\T\\B\n & & NLO+NNLL${}^{\\dagger}$ & $15.74 \\pm 0.40~\\text{(scale)} {}^{+1.10}_{-1.14}~\\text{(PDF+$\\alpha_{\\rm s}$)} \\pm 0.28~\\text{($m_\\ensuremath{t}\\xspace$)} \\pm 0.04~\\text{($E_{\\rm beam}$)}$\\\\\n \\cline{3-4}\n \\cline{2-4}\n \\T\\B\n & \\multirow{2}{*}{\\mys-channel\\xspace} & NLO${}^{\\dagger}$ & $\\;\\,4.29_{-0.10}^{+0.12}~\\text{(scale)} \\pm 0.14~\\text{(PDF+$\\alpha_{\\rm s}$)} \\pm 0.10~\\text{($m_\\ensuremath{t}\\xspace$)} \\pm 0.01~\\text{($E_{\\rm beam}$)}$ \\\\\n \\cline{3-4}\n \\T\\B\n & & NLO+NNLL & $\\;\\,4.63^{+0.20}_{-0.18}~\\text{(scale+PDF+$\\alpha_{\\rm s}$)}$\\\\\n \\cline{3-4}\n \\hline\n \\hline\n \\T\\B\n & & NLO${}^{\\dagger}$ & $84.7_{-1.7}^{+2.6}~\\text{(scale)} \\pm 2.8~\\text{(PDF+$\\alpha_{\\rm s}$)} \\pm 0.8~\\text{($m_\\ensuremath{t}\\xspace$)} \\pm 0.2~\\text{($E_{\\rm beam}$)}$ \\\\\n \\cline{3-4}\n \\T\\B\n & \\myt-channel\\xspace & NLO+NNLL & $87.8^{+3.4}_{-1.9}~\\text{(scale+PDF+$\\alpha_{\\rm s}$)}$\\\\\n \\cline{3-4}\n \\T\\B\n & & NNLO & $84.2^{+0.3}_{-0.2}~\\text{(scale)}$\\\\\n \\cline{2-4}\n \\T\\B\n \\ensuremath{ 8~\\TeV} & \\multirow{2}{*}{\\ensuremath{tW}\\xspace} & NLO & $18.77_{-0.82}^{+0.77}~\\text{(scale)} \\pm 1.70~\\text{(PDF+$\\alpha_{\\rm s}$)}$ \\\\\n \\cline{3-4}\n \\T\\B\n & & NLO+NNLL${}^{\\dagger}$ & $22.37 \\pm 0.60~\\text{(scale)} \\pm1.40~\\text{(PDF+$\\alpha_{\\rm s}$)} \\pm 0.38~\\text{($m_\\ensuremath{t}\\xspace$)} \\pm 0.06~\\text{($E_{\\rm beam}$)}$\\\\\n \\cline{3-4}\n \\cline{2-4}\n \\T\\B\n & \\multirow{2}{*}{\\mys-channel\\xspace} & NLO${}^{\\dagger}$ & $\\;\\,5.24_{-0.12}^{+0.15}~\\text{(scale)} \\pm 0.16~\\text{(PDF+$\\alpha_{\\rm s}$)} \\pm 0.12~\\text{($m_\\ensuremath{t}\\xspace$)} \\pm 0.01~\\text{($E_{\\rm beam}$)}$ \\\\\n \\cline{3-4}\n \\T\\B\n & & NLO+NNLL & $\\;\\,5.61\\pm0.22~\\text{(scale+PDF+$\\alpha_{\\rm s}$)}$\\\\\n \\cline{3-4}\n \\hline\n \\hline\n \\end{tabular}\n \\IfPackagesLoaded{adjustbox}{\\end{adjustbox}}{}\n \\label{tab:xs}\n\\end{table}\n\nIn this paper, NLO predictions serve as the reference for the $t$- and $s$-channel processes, following the prescriptions presented above, because higher-order calculations and their uncertainties are not fully computed and available for the parameter values of choice. The advantage of the NLO cross-section calculations is that the configurable parameters in \\textsc{HatHor}\\xspace can be set according to those used to generate the ATLAS and CMS simulation samples. The $t$- and \\mys-channel\\xspace\\ processes do not interfere at NLO~\\cite{Willenbrock:1986cr}. For these two processes, the entire phase space is included in the integration in order to obtain the total cross-section. The \\ensuremath{tW}\\xspace cross-section prediction, $\\ensuremath{\\sigma^{\\mytW}_{\\mathrm{\\textrm{theo.}}}}\\xspace$, is available at NLO ~\\cite{Kant:2014oha} and NLO+NNLL~\\cite{Kidonakis:2013zqa,Kidonakis:2010ux}. The \\ensuremath{tW}\\xspace process at NLO interferes with the \\ensuremath{t \\bar{t}}\\xspace process at LO with the subsequent decay $\\ensuremath{\\bar{t}}\\xspace \\rightarrow W\\ensuremath{\\bar{b}}\\xspace$. In the NLO prediction for \\ensuremath{tW}\\xspace production provided in Ref.~\\cite{Kant:2014oha}, a kinematic cut-off is imposed on the transverse momentum (\\pt) of the outgoing \\bquark, suppressing the contribution from \\ttbar\\ production. Since the treatment of this interference in \\textsc{HatHor}\\xspace is still being developed~\\cite{Demartin:2016axk, Campbell:2005bb}, the NLO+NNLL calculation is used as reference for \\ensuremath{tW}\\xspace production. For the reference cross-section predictions, uncertainties corresponding to the dependence on $m_{\\myt}\\xspace$ and on the LHC beam energy, $E_{\\rm beam}$, are evaluated. The $m_{\\myt}\\xspace$ dependence is estimated by varying its central value of $172.5~\\GeV$ (the value used in the simulation samples used to measure the single-top-quark cross-sections) by $\\pm$1~\\GeV, using the functional form proposed in Ref.~\\cite{Czakon:2013goa}. The theoretical calculations are performed at a given centre-of-mass energy while the energy of the LHC beam is measured with an uncertainty. The single-top-quark cross-sections are assumed to depend on $E_{\\rm beam}$ according to the model given in Ref.~\\cite{Langenfeld:2009tc}, with a relative uncertainty $\\delta E_{\\rm beam}\/E_{\\rm beam}$ of 0.1\\%~\\cite{Todesco:2017nnk}. \nThe theoretical cross-sections that are used as reference are marked with a $^\\dagger$ in Table~\\ref{tab:xs}.\n\n\n\n \\section{Single-top-quark cross-section measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}}\n\\label{sec:atlascmsxs}\n\nThe \\myt-channel\\xspace single-top-quark production cross-sections, $\\ensuremath{\\sigma_{{\\myt\\textrm{-chan.}}}}\\xspace$, were measured by the\nATLAS and CMS Collaborations at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}~\\cite{TOPQ-2012-21,CMS-TOP-11-021}\nand 8\\,\\TeV~\\cite{TOPQ-2015-05,CMS-TOP-12-038}. Evidence of \\ensuremath{tW}\\xspace production was reported at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}~by \nATLAS~\\cite{TOPQ-2011-17} and CMS~\\cite{CMS-TOP-11-022}, \nwhile at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ its cross-section, $\\ensuremath{\\sigma_{\\mytW}}\\xspace$, was\nmeasured by both\nexperiments~\\cite{TOPQ-2012-20,CMS-TOP-12-040}. Evidence of \\mys-channel\\xspace\\\nproduction was reported by ATLAS, with a measured cross-section,\n$\\ensuremath{\\sigma_{\\mys\\textrm{-chan.}}}\\xspace$, at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}~\\cite{TOPQ-2015-01}, whereas CMS set upper limits on\nthe \\mys-channel\\xspace\\ production cross-section at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}. The observed\n(expected) significance of the CMS measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ is $2.3$ ($0.8$) standard\ndeviations~\\cite{CMS-TOP-13-009}.\n\nThe ATLAS and CMS analyses use similar approaches to measure the single-top-quark production cross-sections. Both experiments select events containing at least one prompt isolated lepton (electron or muon) and at least one high-\\pt\\ jet. The analyses use various multivariate analysis (MVA) techniques, such as boosted decision trees~\\cite{FREUND1997119,Friedman2006,Hocker:2007ht}, neural networks~\\cite{Feindt:2006pm}, or the matrix element method (MEM)~\\cite{Kondo:1988yd, Kondo:1991dw}, to separate the signal from background. To measure the cross-section, analyses perform a binned maximum-likelihood fit to data using the distribution of the corresponding MVA discriminator. \nExceptions are the ATLAS \\mys-channel\\xspace\\ and CMS \\myt-channel\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}. In the ATLAS \\mys-channel\\xspace\\ analysis, the fit is performed simultaneously to the MEM discriminant in the signal region and the lepton-charge distribution in the \\ensuremath{W}+jets\\xspace control region. The CMS \\myt-channel\\xspace\\ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ is based on a simultaneous fit to the absolute pseudorapidity ($\\eta$) distributions of the recoiling light-flavour jet in events with negative and with positive lepton charge. \nThe analyses measuring different single-top-quark production modes within the same experiment and at the same centre-of-mass energy have disjoint signal regions. \nBoth experiments simulate the single-top-quark processes using the NLO \\POWHEGBOX\\ generator~\\cite{Nason:2004rx,Frixione:2007vw,Alioli:2009je,Alioli:2010xd,Re:2010bp} for the matrix-element (ME) calculations. ATLAS also uses the \\POWHEGBOX\\ generator to simulate top-quark-pair background events, while CMS uses the LO \\MADGRAPH generator~\\cite{Alwall:2011uj}. The \\PYTHIA~\\cite{Sjostrand:2006za} event generator is used for modelling the parton shower (PS), hadronisation and the underlying event in both the single-top-quark and \\ttbar\\ processes. The cross-sections are measured assuming a value of $172.5$~\\GeV\\ for $m_{\\myt}\\xspace$ for all top-quark processes\nand all centre-of-mass energies. \nA summary of the uncertainties in each measurement is shown in Table~\\ref{tab:xsec_meas}, with details given in Appendix~\\ref{app:allUncs}.\n\n\\begin{table}[!htbp]\n \\centering\n \\caption{Summary of the single-top-quark cross-section measurements published by the ATLAS and\n CMS Collaborations at\n \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}. Total uncertainties are shown. Small\n differences between the integrated luminosity \n values in different analyses within the same experiment and\n centre-of-mass energy are due to different luminosity\n calibrations at the time of publication.}\n \\IfPackagesLoaded{adjustbox}{\\begin{adjustbox}{max width=1.0\\textwidth}}{}\n \\begin{tabular}{l|l|r@{\\hskip 0cm}lc|r@{\\hskip 0cm}lc}\n \\hline\n \\hline\n \\multicolumn{2}{c|}{} & \\multicolumn{3}{c|}{ATLAS} & \\multicolumn{3}{c}{CMS} \\\\\n \\hline\n $\\sqrt{s}$ & Process & \\multicolumn{2}{c}{$\\sigma{}$\\,[pb]} & Lumi.\\ [$\\ensuremath{\\rm fb^{-1}}\\xspace$] & \\multicolumn{2}{c}{$\\sigma{}$\\,[pb]} & Lumi.\\ [$\\ensuremath{\\rm fb^{-1}}\\xspace$] \\\\\n \\hline\n & \\myt-channel\\xspace & $68$&$~\\pm~8$ & 4.59 & $67.2$&$~\\pm~6.1$ & 1.17--1.56 \\\\ \n \\cline{2-8}\n \\T\\B\n \\ensuremath{ 7~\\TeV} & \\ensuremath{tW}\\xspace & $16.8$&$~\\pm~5.7$ & 2.05 & \\multicolumn{2}{c}{$16^{+5}_{-4}$} & 4.9 \\\\\n \\cline{2-8}\n \\T\\B\n & \\mys-channel\\xspace & \\multicolumn{2}{c}{---} & --- & $7.1$&$~\\pm~8.1$ & 5.1 \\\\\n \\hline\n \\hline\n \\T\\B\n & \\myt-channel\\xspace & \\multicolumn{2}{c}{$89.6^{+7.1}_{-6.3}$} & 20.2 & $83.6$&$~\\pm~7.8$ & 19.7 \\\\\n \\cline{2-8}\n \\T\\B\n \\ensuremath{ 8~\\TeV} & \\ensuremath{tW}\\xspace & \\multicolumn{2}{c}{$23.0^{+3.6}_{-3.9}$} & 20.3 & $23.4$&$~\\pm~5.4$ & 12.2 \\\\\n \\cline{2-8}\n \\T\\B\n & \\mys-channel\\xspace & \\multicolumn{2}{c}{$4.8^{+1.8}_{-1.5}$} & 20.3 & $13.4$&$~\\pm~7.3$ & 19.7 \\\\\n \\hline\n \\hline\n \\end{tabular}\n \\renewcommand{\\arraystretch}{1}\n \\IfPackagesLoaded{adjustbox}{\\end{adjustbox}}{}\n \\label{tab:xsec_meas}\n\\end{table}\n \\section{Combination methodology}\n\\label{sec:method}\nThe ATLAS and CMS single-top-quark production cross-section measurements shown in Table~\\ref{tab:xsec_meas} are combined, and the combined $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ value determined, using the best linear unbiased estimator (BLUE) method~\\cite{Lyons:1988rp, Valassi:2003mu,Nisius:2014wua}. The BLUE method is applied iteratively in order to reduce a possible bias arising from the dependence of systematic uncertainties on the central value of the cross-section~\\cite{Lista:2014qia}. Convergence is reached when the central value changes by less than 0.01\\% compared with the previous iteration. \nIn each iteration, the BLUE method minimises the global $\\chi^2$ by adjusting the weight for each input measurement~\\cite{,Nisius:2014wua}. The global $\\chi^2$ is calculated taking correlations\ninto account. The sum of weights is required to be equal to one. Negative weights are allowed; these indicate strong correlations~\\cite{Valassi:2013bga}. The number of degrees of freedom is $n-1$, where $n$ is number of measurements in the combination. The $\\chi^2$ and $n$ are then used to calculate a corresponding probability~\\cite{Nisius:2014wua}. The systematic uncertainties are scaled with the cross-section in each iteration, i.e.\\ they are treated as relative uncertainties. The data and simulation statistical uncertainties are not scaled~\\cite{Lista:2014qia}. The systematic uncertainties in the \\mys-channel\\xspace\\ cross-section combination are also not scaled because the \\mys-channel\\xspace\\ measurements have large backgrounds.\n\nFollowing the same strategy as in the input measurements by the ATLAS and CMS Collaborations, the combined cross-sections are reported at $m_{\\myt}\\xspace = 172.5~\\GeV$, not including the uncertainty associated with the $m_{\\myt}\\xspace$ variation. The shift in the combined cross-section due to a variation of $\\pm 1~\\GeV$ in the top-quark mass is given where this information is available. For the determination of the combined $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ value, the uncertainty in the measured cross-sections due to a variation of $\\pm 1~\\GeV$ in the mass is considered.\nUncertainties in the measurements are symmetrised, before combination, by averaging the magnitude of the downward and upward variations. More details are given in Sections~\\ref{sec:systcat} and \\ref{sec:xscomb}.\n \\section{Systematic uncertainties and correlation assumptions}\n\\label{sec:systcat}\nIn order to combine single-top-quark cross-section measurements and $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ values, the sources of uncertainty are grouped into categories. While the categorisation and evaluation of uncertainties varies somewhat between experiments and between measurements, each individual measurement considers a complete set of uncertainties. Assumptions are made about correlations between similar sources of uncertainty in different measurements, as explained in Section~\\ref{sec:systcat_exp}. Uncertainties associated with theoretical predictions are taken into account in the $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ combination. The correlations between similar uncertainties in different theoretical predictions are discussed in Section~\\ref{sec:systcat_theo}.\n\n\\subsection{Systematic uncertainties in measured cross-sections}\n\\label{sec:systcat_exp}\n\nSystematic uncertainties in the ATLAS \\myt-channel\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ are evaluated using pseudoexperiments, except the background normalisation uncertainties, which are constrained in the fit to data. In the ATLAS \\ensuremath{tW}\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ and the \\mys-channel\\xspace\\ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, systematic uncertainties are included as nuisance parameters in profile-likelihood fits. Systematic uncertainties in the CMS \\myt-channel\\xspace\\ and \\ensuremath{tW}\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ are included as nuisance parameters in fits to data, except the theory modelling uncertainties in signal and backgrounds, described below, which are evaluated using pseudoexperiments. All systematic uncertainties in the CMS \\mys-channel\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ are obtained through pseudoexperiments, except the background normalisation uncertainties, which are constrained in the fit to data. In the analyses where systematic uncertainties are included as nuisance parameters, the total uncertainty presented in Table~\\ref{tab:xsec_meas} is evaluated by varying all the nuisance parameters in the fit simultaneously. To extract the impact of each source of this type of uncertainty, these analyses use approximate procedures which neglect the correlations between sources of uncertainty introduced by the fits. Throughout this paper, individual uncertainties are taken as reported by the input analyses, regardless of the method used to determine them. The total uncertainties are evaluated as the sum in quadrature of individual contributions.\n\nAlthough the sources of systematic uncertainty and the procedures used to\nestimate their impact on the measured cross-section are partially different in the individual analyses, it is still possible to identify contributions that describe similar physical effects. These contributions are listed below; they are grouped together, and only the resulting categories are used in the combination. Categories are treated as uncorrelated among each other. \nFor each source of uncertainty, correlations between different measurements are assumed to be positive, unless explicitly mentioned otherwise. \nThe stability of the cross-section and $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ combinations is studied by varying the correlation assumptions for the dominant uncertainties, as discussed in Section~\\ref{sec:tests}.\n\nThe uncertainties in each category are listed below, with the correlation assumptions across experiments given in parentheses. These correlations correspond to those used in the cross-section combinations. They are also valid for the combination of the $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ extractions, unless explicitly mentioned otherwise. The symbol ``---'' means that the uncertainty is either considered only in the ATLAS or the CMS measurement, or is not considered at all. A summary of uncertainties in the cross-section measurements together with the corresponding correlation assumptions between experiments is provided in Appendix~\\ref{app:allUncs}.\n\n\\noindent {\\bf Data statistical} (Correlation 0)\\\\ This statistical uncertainty arises from the limited size of the data sample. It is uncorrelated between ATLAS and CMS, between production modes, and between centre-of-mass energies.\n\n\\noindent {\\bf Simulation statistical} (Correlation 0 and --- for CMS \\ensuremath{tW}\\xspace\\ at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ and \\mys-channel\\xspace\\ at \\ensuremath{\\sqrt{s}=\\lhcenergyEight})\\\\ This statistical uncertainty comes from the limited size of simulated event samples. It is uncorrelated between ATLAS and CMS, between production modes, and between centre-of-mass energies. For the CMS \\ensuremath{tW}\\xspace\\ analysis at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ and \\mys-channel\\xspace\\ analysis at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, this uncertainty is evaluated as part of the total statistical uncertainty, which is also considered uncorrelated, as discussed above. \nMore details are given in Appendices~\\ref{app:tw} and \\ref{app:sch}.\n\n\\noindent {\\bf Integrated luminosity} (Correlation 0.3)\\\\ This uncertainty originates from the systematic uncertainty in the integrated luminosity, as determined by the individual experiments using the methods described in Refs.~\\cite{DAPR-2011-01,DAPR-2013-01,CMS-PAS-SMP-12-008,CMS-PAS-LUM-13-001}. It affects the determination of both the signal and background yields. The integrated-luminosity uncertainty has a component that is correlated between ATLAS and CMS, arising from imperfect knowledge of the beam currents during van der Meer scans in the LHC accelerator~\\cite{vanderMeer:1968zz}, and an uncorrelated component from the long-term luminosity monitoring that is experiment-specific.\nAt \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, these components are 0.5\\% and 1.7\\% respectively for ATLAS and 0.5\\% and 2.1\\% respectively for CMS. At \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, they are 0.6\\% and 1.8\\% respectively for ATLAS and 0.7\\% and 2.5\\% respectively for CMS.\nAt both centre-of-mass energies, the correlation coefficient between the integrated-luminosity uncertainty in ATLAS and CMS at the same centre-of-mass energy is $\\rho=0.3$. Within the same experiment, the integrated-luminosity uncertainty is assumed to be correlated between production modes and uncorrelated between centre-of-mass energies. In Section~\\ref{sec:tests}, it is shown that the combined $\\ensuremath{\\absvtb^2}\\xspace$ result does not depend significantly on the correlation assumptions.\n\n\\noindent {\\bf Theory modelling}\\\\ This category contains the uncertainties in the modelling of the simulated single-top-quark processes, as well as smaller contributions from the modelling of the \\ensuremath{t \\bar{t}}\\xspace and \\ensuremath{W}+jets\\xspace background processes. Both signal and background modelling are included because the uncertainties in all top-quark processes are closely related.\n These include initial- and final-state radiation (ISR\/FSR), renormalisation and factorisation scales, NLO matching method, PS and hadronisation modelling, and PDF uncertainties. For the \\ensuremath{tW}\\xspace process, the uncertainty due to the treatment of interference between \\ensuremath{tW}\\xspace and \\ensuremath{t \\bar{t}}\\xspace final states is also included, as discussed below. These modelling uncertainties in signal and background processes are summed in quadrature in each input measurement.\n\n\\begin{itemize}\n\n\\item {\\it Scales and radiation modelling} (Correlation 1)\\\\ The renormalisation and factorisation scales and ISR\/FSR uncertainties account for missing higher-order corrections in the perturbative expansion and the amount of initial- and final-state radiation in simulated signal and background processes. In the ATLAS measurements of all three production modes, these uncertainties are estimated using dedicated single-top-quark and \\ensuremath{t \\bar{t}}\\xspace simulated event samples, by consistently varying the renormalisation and factorisation scales and the amount of ISR\/FSR in accordance with a measurement of additional jet activity in \\ensuremath{t \\bar{t}}\\xspace events at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}~\\cite{TOPQ-2011-21,ATL-PHYS-PUB-2014-005}.\nIn the ATLAS \\myt-channel\\xspace\\ measurements, they are also estimated in \\ensuremath{W}+jets\\xspace\\ simulated event samples, by varying the scale and matching parameters in the \\ALPGEN\\ LO multileg generator~\\cite{Mangano:2002ea} at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ and by varying the parameters controlling the scale in the \\SHERPA\\ LO multileg generator~\\cite{Gleisberg:2008ta} at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}. In the CMS measurements, these uncertainties are estimated by varying the renormalisation and factorisation scales, and ISR\/FSR, consistently in the simulated event samples. In the CMS \\myt-channel\\xspace\\ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, this uncertainty applies only to the signal modelling since the modelling of the dominant \\ensuremath{t \\bar{t}}\\xspace and \\ensuremath{W}+jets\\xspace background processes is obtained from data. However, for the \\myt-channel\\xspace\\ analysis at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, the scales are varied in the simulated signal, \\ensuremath{t \\bar{t}}\\xspace, \\ensuremath{W}+jets\\xspace and other single-top-quark processes. The same approach is followed in the CMS \\mys-channel\\xspace\\ measurements at both centre-of-mass energies. The \\ensuremath{tW}\\xspace cross-section measurements of CMS account for this uncertainty only in the \\ensuremath{tW}\\xspace signal and \\ensuremath{t \\bar{t}}\\xspace background, given the negligible contributions from the \\ensuremath{W}+jets\\xspace and other single-top-quark processes in the dilepton final state.\n\nAlthough the methods are apparently different, they mostly address the same uncertainty, hence this uncertainty is considered correlated between ATLAS and CMS. It is also considered correlated between production modes and centre-of-mass energies. The combined $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ result does not depend significantly on this correlation assumption,\nas discussed in Section~\\ref{sec:tests}.\n\n\n\\item {\\it NLO matching} (Correlation 1 for \\myt-channel\\xspace\\ and --- for \\ensuremath{tW}\\xspace and \\mys-channel\\xspace)\\\\ The ATLAS measurements include an uncertainty to account for different NLO matching methods implemented in different NLO event generators. This is evaluated in single-top-quark and \\ensuremath{t \\bar{t}}\\xspace simulations by comparing the \\POWHEGBOX, \\MCatNLO~\\cite{Frixione:2002ik,Skands:2010ak}, and \\MGMCatNLO~\\cite{Alwall:2014hca} event generators, all interfaced to \\HERWIG~\\cite{Corcella:2000bw} (with \\JIMMY~\\cite{Butterworth:1996zw} for the underlying-event modelling). In the CMS \\myt-channel\\xspace\\ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, the NLO matching uncertainty is evaluated by comparing \\POWHEGBOX\\ with \\Comphep~\\cite{Boos:2004kh,Pukhov:1999gg}. In the CMS \\myt-channel\\xspace\\ analysis at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, this uncertainty accounts for different NLO matching methods in the \\myt-channel\\xspace\\ signal event generator, as well as for differences between event generation in the 4FS and 5FS, by comparing \\POWHEGBOX\\ with \\MADGRAPH. The NLO matching uncertainty is considered correlated between ATLAS and CMS, between production modes, and between centre-of-mass energies. In the CMS \\ensuremath{tW}\\xspace and \\mys-channel\\xspace\\ analyses at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}, this uncertainty is not considered, since the modelling uncertainties in the scheme to remove overlap with \\ensuremath{t \\bar{t}}\\xspace are dominant in the \\ensuremath{tW}\\xspace analysis and the renormalisation\/factorisation scale is dominant in the \\mys-channel\\xspace\\ analysis. The results of the stability test for this uncertainty are shown in Section~\\ref{sec:tests}.\n\n\\item {\\it Parton shower and hadronisation} (Correlation 1)\\\\ In both experiments, the difference between the \\PYTHIA and \\HERWIG showering programs is considered in the jet energy scale (JES)~\\cite{PERF-2011-03,PERF-2012-01,CMS-JME-10-011,CMS-JME-13-004} and \\btag calibration~\\cite{PERF-2012-04,ATLAS-CONF-2014-004,ATLAS-CONF-2014-046,CMS-BTV-12-001,CMS-PAS-BTV-13-001}. \nThe ATLAS analyses additionally include an uncertainty in the PS and hadronisation modelling in simulated single-top-quark and \\ensuremath{t \\bar{t}}\\xspace events, evaluated by comparing the \\POWHEGBOX\\ event generator interfaced to \\PYTHIA or to \\HERWIG. \nThe CMS analyses additionally include an uncertainty in the \\ensuremath{t \\bar{t}}\\xspace and \\ensuremath{W}+jets\\xspace backgrounds estimated with the \\MADGRAPH event generator interfaced to \\PYTHIA. It is evaluated in simulated event samples where the value of the ME\/PS matching threshold in the MLM method~\\cite{Mangano:2006rw} is doubled or halved from its initial value. \nThe CMS \\myt-channel\\xspace\\ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ does not consider this uncertainty in the \\ensuremath{t \\bar{t}}\\xspace and \\ensuremath{W}+jets\\xspace backgrounds since the distribution and normalisation of the \\ensuremath{t \\bar{t}}\\xspace\\ and \\ensuremath{W}+jets\\xspace\\ processes are derived mostly from data.\nIn the CMS \\ensuremath{tW}\\xspace analyses at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}, the contributions of the \\ensuremath{W}+jets\\xspace and other single-top-quark processes in the dilepton final state are negligible.\n\nThis uncertainty is considered correlated between ATLAS and CMS, between different production modes, and between different centre-of-mass energies. The combined $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ result does not depend significantly on this correlation assumption, as shown in Section~\\ref{sec:tests}.\n\n\\item {\\it Parton distribution functions} (Correlation 1)\\\\ The PDF uncertainty is evaluated following the PDF4LHC procedures~\\cite{Alekhin:2011sk,Botje:2011sn,Butterworth:2015oua} and is \nconsidered correlated between ATLAS and CMS, between different production modes, and between different centre-of-mass energies.\n\n\\item {\\it \\ensuremath{tW}\\xspace and \\ensuremath{t \\bar{t}}\\xspace interference} (Correlation 1 for \\ensuremath{tW}\\xspace\\ and --- for $t$- and $s$-channels)\\\\ The \\ensuremath{tW}\\xspace process interferes with \\ensuremath{t \\bar{t}}\\xspace production at NLO~\\cite{Frixione:2008yi, tWth2, tWth3}. In both ATLAS and CMS, two simulation approaches are compared: diagram removal (DR)~\\cite{Frixione:2008yi} and diagram subtraction (DS)~\\cite{Frixione:2008yi,Tait:1999cf}.\nIn the DR approach, all NLO diagrams that overlap with the doubly resonant \\ensuremath{t \\bar{t}}\\xspace\\ contributions are removed from the calculation of the \\ensuremath{tW}\\xspace\\ amplitude. This approach accounts for the interference term, but it is not gauge invariant (though the effect is numerically negligible)~\\cite{Frixione:2008yi}. In the DS approach, a subtraction term is built into the amplitude to cancel out the \\ensuremath{t \\bar{t}}\\xspace\\ component close to the top-quark resonance while respecting gauge invariance.\n\n\n The DR approach is the default, and the comparison with the DS approach is used to assess this systematic uncertainty. For the \\ensuremath{tW}\\xspace analyses, this uncertainty is considered correlated between the two experiments and between different centre-of-mass energies.\n\n\\item {\\it Modelling of the top-quark \\ensuremath{p_{\\mathrm{T}}}\\xspace\\ spectrum} (Correlation ---)\\\\ In the CMS \\ensuremath{tW}\\xspace and \\mys-channel\\xspace\\ analyses at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, the simulated \\ensuremath{t \\bar{t}}\\xspace events are reweighted to correct the $\\ensuremath{p_{\\mathrm{T}}}\\xspace$\\ spectrum of the generated top quarks, which was found to be significantly harder than the spectrum observed in data in differential cross-section measurements~\\cite{CMS-TOP-12-028,Czakon:2016ckf}. To estimate the uncertainty related to this mismodelling, the \\ensuremath{tW}\\xspace\\ measurement is repeated without the reweighting, and the change relative to the default result is taken as the uncertainty. In the CMS \\mys-channel\\xspace\\ analysis, the measurement is repeated with the effect of the weights removed and doubled. The resulting variation in the cross-section is symmetrised. This uncertainty is not considered in the CMS \\myt-channel\\xspace\\ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ where the modelling of the \\ensuremath{t \\bar{t}}\\xspace\\ background is extracted from data. In the ATLAS measurements, modelling uncertainties in the top-quark \\ensuremath{p_{\\mathrm{T}}}\\xspace spectrum in \\ensuremath{t \\bar{t}}\\xspace events~\\cite{TOPQ-2015-06} are covered by the PS and hadronisation uncertainty and they are found to be small in comparison with other systematic uncertainties. This uncertainty is considered correlated between the CMS \\ensuremath{tW}\\xspace and \\mys-channel\\xspace\\ analyses at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}.\n\n\\item {\\it Dependence on the top-quark mass} (Correlation 1)\\\\ The measured single-top-quark cross-sections shown in Table~\\ref{tab:xsec_meas} assume a nominal $m_{\\myt}\\xspace$ value of 172.5~\\GeV. The dependence of the measured cross-section on $m_{\\myt}\\xspace$ is estimated for the ATLAS \\myt-channel\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ and for the ATLAS \\ensuremath{tW}\\xspace\\ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}. It is determined using dedicated simulations of single-top-quark and \\ttbar\\ samples with different $m_t$ values. The cross-section measurements assuming the different $m_t$ values are interpolated using a first- or a second-order polynomial, for which the constant term is given by the central value of $m_{\\myt}\\xspace = 172.5$~\\GeV. \n The CMS measurements at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ provide information for a variation of $\\pm 2~\\GeV$ in the top-quark mass, which is scaled to a $\\pm 1~\\GeV$ shift assuming a linear dependence. For the CMS \\myt-channel\\xspace\\ and \\ensuremath{tW}\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, the changes in cross-sections are symmetrised and reported as uncertainties. In the CMS \\mys-channel\\xspace\\ analysis, the change in the cross-section is determined for the up and down variation of $m_{\\myt}\\xspace$.\n No estimates are available for the CMS \\myt-channel\\xspace\\ analysis at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, the ATLAS and CMS \\ensuremath{tW}\\xspace\\ analyses at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ or the ATLAS \\mys-channel\\xspace\\ analysis at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}. The top-quark-mass uncertainty is small for each measurement, thus the impact of not evaluating it for these measurements is negligible.\n\n In this paper, a symmetrised uncertainty in the measured cross-section due to a variation of $\\pm 1~\\GeV$ in the top-quark mass is considered. When the full cross-section dependence on the top-quark mass is available for a given production mode at a given centre-of-mass energy, the sign of the dependence of the uncertainty per unit of mass is taken into account in the correlations.\nIn the case of the CMS \\myt-channel\\xspace\\ and \\ensuremath{tW}\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, where the sign of the dependence is not available, it is assumed that the sign is the same as for the ATLAS measurement, since the phase space and background composition are comparable between CMS and ATLAS. Given that the uncertainty in the measured cross-section is considered for the same $m_t$ variation and considering the sign of the dependence when available, this uncertainty is considered correlated between ATLAS and CMS and between different centre-of-mass energies and uncorrelated between the \\myt-channel\\xspace\\ and \\ensuremath{tW}\\xspace\\ production modes. \n\n\\end{itemize}\n\n\\noindent {\\bf Background normalisation} (Correlation 0)\\\\\nThree background uncertainties are considered: in top-quark background (\\ttbar\\ and other single-top-quark processes), in other background determined from simulation (\\ensuremath{W\/Z}+jets\\xspace, diboson, and other smaller background channels), and in background estimated from data (multijet background from misidentified and non-prompt leptons). The exceptions are the \\myt-channel\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, where the background from simulation includes top-quark background, as shown in Tables~\\ref{tab:tChan7XS}$-$\\ref{tab:sChan8XS} in Appendix~\\ref{app:allUncs}. \nThe normalisation of the main background processes is determined from data, either by inclusion of normalisation uncertainties as nuisance parameters in the fit used to extract the signal, or through dedicated techniques based on data. In the \\myt-channel\\xspace\\ and \\mys-channel\\xspace\\ measurements, the uncertainties in the theoretical cross-section predictions for the top-quark, \\ensuremath{W\/Z}+jets\\xspace, and diboson processes are included. In the \\ensuremath{tW}\\xspace\\ measurements, the uncertainties in the theoretical cross-section predictions for the top-quark and diboson processes are taken into account. In the ATLAS measurements of the \\myt-channel\\xspace\\ process at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}, the uncertainty in the multijet background is estimated by comparing background estimates made using different techniques based on simulation and data samples.\nIn the ATLAS \\ensuremath{tW}\\xspace\\ analyses at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}, the normalisation uncertainty in the background from misidentified and non-prompt leptons is obtained from variations in the data-based estimate.\nIn the ATLAS \\mys-channel\\xspace\\ analysis, the uncertainty assigned to the normalisation of the multijet background is based on control samples.\nFor all CMS measurements, background normalisations are constrained in the fits to data. In the CMS measurements of the \\myt-channel\\xspace\\ and \\mys-channel\\xspace\\ processes, the uncertainties in the multijet background are assessed by comparing the results of alternative background estimation methods based on data. \nHence, the associated uncertainties are considered uncorrelated between ATLAS and CMS, between different production modes, and between different centre-of-mass energies. \n\n\n\\noindent {\\bf Jets} \\\\ In the analyses, the uncertainties related to the reconstruction and energy calibration of jets are propagated through variations in the modelling of the detector response. These uncertainties, classified in categories as JES, jet identification (JetID), and jet energy resolution (JER), are discussed below.\n\n\\begin{itemize}\n\\item {\\it Jet energy scale} (Correlation 0 and --- for JES flavour)\\\\ \nThe JES is derived using information from data and simulation. Its uncertainty increases with increasing $|\\eta|$ and decreases with increasing $\\pT$ of the reconstructed jet.\n\nFor all of the ATLAS measurements, except the $\\ensuremath{tW}\\xspace$ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, the JES uncertainty is split into components originating from the jet calibration procedure; most of them are derived from in situ techniques based on data~\\cite{PERF-2011-03,PERF-2012-01}. These components are categorised as modelling, detector, calibration method, and statistical components, which are grouped into the ``JES common'' uncertainty, as well as a flavour-dependence component (``JES flavour''), which accounts for the flavour composition of the jets and the calorimeter response to jets of different flavours. The modelling of additional $pp$ collisions in each bunch-crossing (pile-up) is considered separately, as discussed below. The $\\eta$-dependent component is dominant for the \\myt-channel\\xspace\\ production mode. Thus, the JES common uncertainty is considered uncorrelated between the \\myt-channel\\xspace\\ and the other single-top-quark production modes. For the \\ensuremath{tW}\\xspace analysis at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, the modelling component, which is constrained in the fit to data, is dominant. The uncertainty in the flavour composition of the jets is dominant for the \\mys-channel\\xspace. \n\nFor the CMS measurements, sources contributing to the JES uncertainty are combined together into the ``JES common'' uncertainty, and the effect is propagated to the cross-section measurements through $\\eta$- and $\\ensuremath{p_{\\mathrm{T}}}\\xspace$-dependent JES uncertainties~\\cite{CMS-JME-10-011,CMS-JME-13-004}. The jet energy corrections and their corresponding uncertainties are extracted from data. The JES uncertainty is estimated from its effect on the normalisation and shape of the discriminant in each analysis. The JES uncertainty is considered uncorrelated between the \\myt-channel\\xspace\\ and the other single-top-quark production modes because it is dominated by the forward jet in the \\myt-channel\\xspace.\n\nThe correlation between the JES common uncertainty (or the JES uncertainty for the $\\ensuremath{tW}\\xspace$ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}) in ATLAS and the JES uncertainty in CMS follows the prescription in Refs.~\\cite{ATL-PHYS-PUB-2014-020, ATL-PHYS-PUB-2015-049}, with the slight differences for the \\myt-channel\\xspace\\ described above. The JES common (or JES) uncertainty is considered uncorrelated between ATLAS and CMS, between centre-of-mass energies, and between production modes. Within the ATLAS experiment, the JES common uncertainty is considered correlated between \\ensuremath{tW}\\xspace\\ and \\mys-channel\\xspace\\ and uncorrelated between \\myt-channel\\xspace\\ and the other production modes. For the ATLAS \\myt-channel\\xspace\\ analyses, a correlation of 0.75 is assumed between \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}, since these analyses are mainly affected by the same uncertainty components. This correlation value is estimated by comparing variations of the JES uncertainty components in these two measurements.\n\nIn all CMS measurements and in the ATLAS $\\ensuremath{tW}\\xspace$ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, the JES uncertainty is not split and therefore the JES flavour uncertainty is included in the overall JES uncertainty. For the ATLAS measurements where this component is available, the JES flavour uncertainty is considered correlated between different production modes and uncorrelated between centre-of-mass energies.\n\n\nThe JES uncertainty is one of the dominant contributions in most of the single-top-quark measurements. To ensure the robustness of the results against the correlation assumptions for this large uncertainty, the combination is performed with alternative correlation values, as discussed in Section~\\ref{sec:tests}.\n\n\\item {\\it Jet identification} (Correlation ---)\\\\ In the ATLAS measurements, the JetID uncertainty includes the jet and vertex reconstruction efficiency uncertainties. In the CMS measurements, this uncertainty is included in the JES uncertainty.\n For ATLAS, it is considered correlated between the different production modes at the same centre-of-mass energy and uncorrelated for the other cases.\n \n\\item {\\it Jet energy resolution} (Correlation 0)\\\\ The uncertainty in the JER, which is not split into components, is extracted from data. Generally, the JER uncertainty is propagated via a nuisance parameter in the signal extraction fit, except for the ATLAS \\myt-channel\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}, and the CMS \\mys-channel\\xspace\\ measurement, where this uncertainty is determined using pseudoexperiments. The JER uncertainty is considered uncorrelated between ATLAS and CMS, and between centre-of-mass energies. It is considered correlated between different production modes.\n\\end{itemize}\n\n\\noindent {\\bf Detector modelling} \\\\ This category includes the uncertainty in the modelling of leptons, magnitude of the missing transverse momentum (\\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace), and identification of jets from $b$-quarks (\\btag).\n\n\\begin{itemize}\n\\item {\\it Lepton modelling} (Correlation 0)\\\\ \nThe lepton modelling uncertainty includes components associated with the lepton energy scale and resolution, reconstruction and trigger efficiencies. This uncertainty is considered uncorrelated between ATLAS~\\cite{PERF-2014-05,PERF-2013-05,PERF-2013-03,PERF-2016-01} and CMS~\\cite{Sirunyan:2017ulk} and between different centre-of-mass energies, since it is determined from data. It is considered correlated between different production modes.\n\n\\item {\\it Hadronic part of the high-level trigger} (Correlation ---)\\\\ \nIn the CMS \\myt-channel\\xspace\\ cross-section measurement at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, the high-level trigger (HLT) criteria for the electron channel are based on the presence of an electron together with a \\btagged jet. In this analysis, the uncertainty in the modelling of the hadronic part of the HLT requirement is determined from data. This uncertainty is only evaluated in this one measurement.\n\n\\item {\\it \\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace modelling} (Correlation 0)\\\\ \nThe ATLAS measurements include separate components for the uncertainties in the energy scale and resolution of the \\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace~\\cite{PERF-2014-04}. The CMS measurements account for a combined \\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace scale and resolution uncertainty~\\cite{CMS-JME-10-011,CMS-JME-13-003}, arising from the jet-energy uncertainties.\nAdditionally, CMS accounts for an uncertainty in \\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace arising from energy deposits in the detector that are not included in the reconstruction of leptons, photons, and jets.\nThe \\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace uncertainty is considered uncorrelated between ATLAS and CMS, and between different centre-of-mass energies. It is considered correlated between production modes, except for the ATLAS and CMS \\ensuremath{tW}\\xspace analyses at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, where it is considered uncorrelated with the other production modes because the \\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace uncertainty is constrained in the fit to data. In the ATLAS \\ensuremath{tW}\\xspace analysis at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, this uncertainty is included in the pile-up modelling uncertainty.\n\n\\item {\\it \\btag} (Correlation 0)\\\\ \nIn the ATLAS analyses, \\btag modelling uncertainties are split into components associated with \\bquark, \\ensuremath{c}\\xspace-quark, and light-flavour quark and gluon jets~\\cite{PERF-2012-04,ATLAS-CONF-2014-004,ATLAS-CONF-2014-046}. They are evaluated by varying the $\\pT$-dependence ($\\eta$-dependence in the case of light-flavour jets) of the flavour-dependent scale factors applied to each jet in simulation within a range that reflects the systematic uncertainty in the measured tagging efficiency and misidentification rates. This uncertainty is not considered in the ATLAS \\ensuremath{tW}\\xspace\\ analysis at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ because no \\btag criterion is applied in the event selection. \nIn the CMS measurements, the uncertainties in \\btag efficiency and misidentification rates of jets initiated by light-flavour quarks and gluons are derived from data, using control samples~\\cite{CMS-BTV-12-001,CMS-PAS-BTV-13-001}. The CMS uncertainties are propagated to the cross-section measurements using pseudoexperiments. Exceptions are the \\myt-channel\\xspace\\ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ and the \\ensuremath{tW}\\xspace measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, where these uncertainties are constrained in the fit to data.\n\nThe two collaborations split up the different sources of systematic uncertainties related to \\btag in a different way. However, the different sources are combined by adding their contributions in quadrature to obtain a single \\btag uncertainty per analysis. This means that the \\btag uncertainty also contains the uncertainties associated with the misidentification rates of jets initiated by charm quarks, light-flavour quarks and gluons. The resulting uncertainty is considered uncorrelated between ATLAS and CMS, and between different centre-of-mass energies. It is considered correlated between different production modes.\n\n\\item {\\it Pile-up modelling} (Correlation 0)\\\\ \nIn both ATLAS and CMS, simulated events are reweighted to match the distribution of the average number of interactions per bunch-crossing in data. The corresponding uncertainty is obtained from in situ techniques based on data and simulated event samples. In the ATLAS analyses at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, the uncertainty due to pile-up is derived from the impact of the reweighting on \\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace. In the ATLAS analyses at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, this uncertainty is evaluated as a component of the JES, separated into four terms (number of primary vertices, average number of collisions per bunch-crossing, average pile-up energy density in the calorimeter, and \\pT\\ dependence) since the pile-up calibration (assuming average conditions during 8~\\TeV\\ data-taking) is applied to both data and simulation before selecting and calibrating the jets~\\cite{ATL-PHYS-PUB-2015-049}. \nIn CMS, the reweighting uses a model with a free parameter that can be interpreted as an effective cross-section for inelastic $pp$ interactions. This uncertainty is obtained from a fit to the number of additional primary vertices in simulation. In the CMS analyses, this uncertainty is introduced as a nuisance parameter in the fit. The only exception is the \\mys-channel\\xspace\\ measurement, where the pile-up uncertainty is estimated from pseudoexperiments. In all cases, the effects of pile-up on the jet energy and the isolation of leptons are taken into account in the jet and lepton uncertainties respectively. The pile-up uncertainty is considered uncorrelated between ATLAS and CMS and between different centre-of-mass energies. It is considered correlated between different production modes~\\cite{ATL-PHYS-PUB-2014-020, ATL-PHYS-PUB-2015-049}.\n\n\\end{itemize}\n\n\\subsection{Systematic uncertainties in theoretical cross-section predictions}\n\\label{sec:systcat_theo}\n\nThe systematic uncertainties in the combined $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ value are evaluated from uncertainties in the individual cross-section measurements $\\ensuremath{\\sigma_{\\rm meas.}}\\xspace$ and the theoretical predictions $\\ensuremath{\\sigma_{\\rm theo.}}\\xspace$. The uncertainties associated with $\\ensuremath{\\sigma_{\\rm theo.}}\\xspace$ are discussed in Section~\\ref{sec:theory}; they are summarised in Table~\\ref{tab:xs}. The correlation assumptions for the systematic uncertainties related to the theoretical cross-section are explained below. In Section~\\ref{sec:tests}, the stability of the $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ combination against variations in the correlations is examined. For clarity, the correlations are given in parentheses next to the systematic-uncertainty name. These correlations are used in the combination of the $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ extractions.\\\\\n\n\\noindent {\\bf PDF+$\\alpha_{\\rm s}$} (Correlation 1 for centre-of-mass energies and 0.5 for production modes) \\\\ The PDF uncertainty is considered correlated between centre-of-mass energies and 50\\% correlated between production modes, since different production modes have one initial-state particle in common (a quark or a gluon), but not both.\n\n\\noindent {\\bf Renormalisation and factorisation scales} (Correlation 1 for \\myt-channel\\xspace\\ and \\mys-channel\\xspace\\ and 0 for \\ensuremath{tW}\\xspace) \\\\ The renormalisation and factorisation scale uncertainties in $\\ensuremath{\\sigma_{\\rm theo.}}\\xspace$ are\nconsidered correlated between production modes and centre-of-mass energies, except between the \\ensuremath{tW}\\xspace production mode and the other production modes, where they are considered uncorrelated because the \\ensuremath{tW}\\xspace prediction is computed at a different order in perturbation theory.\n\n\\noindent {\\bf Top-quark mass} (Correlation 1)\\\\ \nThe uncertainty due to $m_{\\myt}\\xspace$ is evaluated by varying $m_{\\myt}\\xspace$ from its central value of 172.5~\\GeV\\ by $\\pm1$~\\GeV\\ and evaluating the corresponding change in cross-section using the parameterisation given in Ref.~\\cite{Czakon:2013goa}, as discussed in Section~\\ref{sec:theory}. This uncertainty is considered correlated between centre-of-mass energies and production modes.\n\n\\noindent {\\bf $E_{\\rm beam}$} (Correlation 1)\\\\ The uncertainty in the cross-section due to the uncertainty in $E_{\\rm beam}$\nis estimated by computing the cross-section variation corresponding to a $\\pm 1$ standard deviation shift in the beam-energy uncertainty.\nIt is considered correlated between centre-of-mass energies and production modes.\n\n\n \\section{Combinations of cross-section measurements}\n\\label{sec:xscomb}\nThe cross-section measurements described in\nSection~\\ref{sec:atlascmsxs} are combined at each centre-of-mass\nenergy for each production mode. Systematic uncertainties are\ncategorised and correlation assumptions are employed according to\nSection~\\ref{sec:systcat}. The combinations are performed using the\niterative BLUE method, as described in Section~\\ref{sec:method}. \n\nAs discussed in Section~\\ref{sec:method}, the uncertainty in the\nmeasured cross-section associated with the $m_{\\myt}\\xspace$ variation is\nnot considered in the combination of cross-sections. However, the\nshift in the combined cross-section resulting from a variation of $\\pm\n1~\\GeV$ in the top-quark mass is provided where this information is\navailable. This is calculated by repeating the combination with the\nup-shifted and down-shifted input cross-sections. In measurements\nwhere only the magnitude of the shift is available for one experiment,\nthe sign of the shift is assumed to be the same for both experiments,\nas discussed in Section~\\ref{sec:systcat_exp}. If the uncertainty associated with the $m_{\\myt}\\xspace$ variation is not\navailable for one or both of the input measurements, then no shift in the combined cross-section is given.\n\n\nAdditional information about the uncertainties considered in the combination of cross-section measurements is provided in Appendix~\\ref{app:allUncs}.\n\n\\subsection{Combinations of \\myt-channel\\xspace\\ cross-section measurements}\n\\label{sec:xscomb_tch}\nThe combination of the ATLAS and CMS \\myt-channel\\xspace\\ cross-section measurements\nat \\ensuremath{\\sqrt{s}=\\lhcenergySeven}~\\cite{TOPQ-2012-21,CMS-TOP-11-021} results, after\none iteration, in \n\\begin{eqnarray*}\n\\ensuremath{\\sigmatch = 67.5\\pm2.4\\;(\\text{stat.})\\pm5.0\\;(\\text{syst.})\\pm1.1\\;(\\text{lumi.})\\,\\,\\text{pb} = 67.5\\pm5.7\\,\\,\\text{pb}}.\n\\end{eqnarray*}\nThe relative uncertainty is $\\ensuremath{8.4\\%}$, which improves\non the uncertainty of 9.1\\% in the most precise individual\nmeasurement from CMS~\\cite{CMS-TOP-11-021}. The $\\chi^2$ for the combination is\n$\\ensuremath{0.01}$, corresponding to a probability of\n$\\ensuremath{93}\\%$. The CMS weight in the combination is\n$\\ensuremath{0.58}$, while the ATLAS weight is\n$\\ensuremath{0.42}$. The overall correlation between the\ntwo measurements is \\ensuremath{20}\\%. The contribution from\neach uncertainty category to the total uncertainty in the combined \\myt-channel\\xspace\\ cross-section \nmeasurement at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ is shown in Table~\\ref{tab:tChanRes}(a).\n\nThe combination of the ATLAS and CMS \\myt-channel\\xspace\\ cross-section measurements\nat \\ensuremath{\\sqrt{s}=\\lhcenergyEight}~\\cite{TOPQ-2015-05,CMS-TOP-12-038} results, after\ntwo iterations, in a cross-section of\n\\begin{eqnarray*}\n\\ensuremath{\\sigmatch = 87.7\\pm1.1\\;(\\text{stat.})\\pm5.5\\;(\\text{syst.})\\pm1.5\\;(\\text{lumi.})\\,\\,\\text{pb} = 87.7\\pm5.8\\,\\,\\text{pb}}.\n\\end{eqnarray*}\nThe relative uncertainty is $\\ensuremath{6.7\\%}$, which improves\non the uncertainty of 7.5\\% in the most precise individual\nmeasurement from ATLAS~\\cite{TOPQ-2015-05}. The $\\chi^2$ for the combination is\n$\\ensuremath{0.59}$, corresponding to a probability of\n$\\ensuremath{44}\\%$. This probability is lower than the\nprobability of the combination at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ because of the differences between the ATLAS and CMS measured cross-sections and their small uncertainties. The ATLAS weight in the combination is\n$\\ensuremath{0.68}$, while the CMS weight is\n$\\ensuremath{0.32}$. The overall correlation between the\ntwo measurements is \\ensuremath{42}\\%. This is larger\nthan the correlation between the measurements at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ because the statistical and detector uncertainties are lower, thus increasing the importance of the theory modelling uncertainty (which is correlated between the two experiments), as shown in Appendix~\\ref{app:tch}. The contribution from each uncertainty category to the\ntotal uncertainty in the combined \\myt-channel\\xspace\\ cross-section measurement at \n\\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ is shown in Table~\\ref{tab:tChanRes}(b).\n\n\\begin{table}[!htbp]\n \\caption{Contribution from each uncertainty category to the combined\n \\myt-channel\\xspace\\ cross-section ($\\ensuremath{\\sigma_{{\\myt\\textrm{-chan.}}}}\\xspace$) uncertainty at\n (a) \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ and (b) \\ensuremath{\\sqrt{s}=\\lhcenergyEight}. The total\n uncertainty is computed by adding in quadrature all the individual systematic\n uncertainties (including the uncertainty in the integrated luminosity) and the statistical uncertainty in data. Correlations of systematic uncertainties between\n experiments are presented in Appendix~\\ref{app:tch}.}\n \\begin{center}\n \\begin{minipage}[b]{0.46\\hsize}\\centering\n \\captionof{subfigure}{}\n \\begin{tabular}{l | r | r}\n\\hline\n \\hline\n \\multicolumn{3}{c}{ $\\ensuremath{\\sigma_{{\\myt\\textrm{-chan.}}}}\\xspace$, $\\sqrt{s} = $ 7~\\TeV } \\\\\n\\hline\n\\hline\n{\\bf Combined cross-section} & \\multicolumn{2}{c}{67.5~pb} \\\\\n\\hline\n\\hline\n\\multirow{2}{*}{Uncertainty category} & \\multicolumn{2}{c}{Uncertainty} \\\\\\cline{2-3}\n & ~~[\\%] & [pb] \\\\\n\\hline\nData statistical & 3.5 & 2.4\\\\\n\\hline\nSimulation statistical & 1.4 & 0.9\\\\\nIntegrated luminosity & 1.7 & 1.1\\\\\nTheory modelling & 5.1 & 3.5\\\\\nBackground normalisation & 1.9 & 1.3\\\\\nJets & 3.4 & 2.3\\\\\nDetector modelling & 3.4 & 2.3\\\\\n\\hline\nTotal syst.\\ unc.\\ (excl.\\ lumi.) & 7.5 & 5.0\\\\\nTotal syst.\\ unc.\\ (incl.\\ lumi.) & 7.6 & 5.2\\\\\n\\hline\\hline\n{\\bf Total uncertainty} & 8.4 & 5.7\\\\\n\\hline\n \\hline\n\\end{tabular}\n \\end{minipage}\n \\hfill\n \\begin{minipage}[b]{0.46\\hsize}\\centering\n \\captionof{subfigure}{}\n \\begin{tabular}{l | r | r}\n\\hline\n \\hline\n \\multicolumn{3}{c}{ $\\ensuremath{\\sigma_{{\\myt\\textrm{-chan.}}}}\\xspace$, $\\sqrt{s} = $ 8~\\TeV } \\\\\n\\hline\n\\hline\n{\\bf Combined cross-section} & \\multicolumn{2}{c}{87.7~pb} \\\\\n\\hline\n\\hline\n\\multirow{2}{*}{Uncertainty category} & \\multicolumn{2}{c}{Uncertainty} \\\\\\cline{2-3}\n & ~~[\\%] & [pb] \\\\\n\\hline\nData statistical & 1.3 & 1.1\\\\\n\\hline\nSimulation statistical & 0.6 & 0.5\\\\\nIntegrated luminosity & 1.7 & 1.5\\\\\nTheory modelling & 5.3 & 4.7\\\\\nBackground normalisation & 1.2 & 1.1\\\\\nJets & 2.6 & 2.3\\\\\nDetector modelling & 1.8 & 1.6\\\\\n\\hline\nTotal syst.\\ unc.\\ (excl.\\ lumi.) & 6.3 & 5.5\\\\\nTotal syst.\\ unc.\\ (incl.\\ lumi.) & 6.5 & 5.7\\\\\n\\hline\\hline\n{\\bf Total uncertainty} & 6.7 & 5.8\\\\\n\\hline\n \\hline\n\\end{tabular}\n \\end{minipage}\n \\end{center}\n \\label{tab:tChanRes}\n\\end{table}\n\nAt both centre-of-mass energies, the uncertainties from theory\nmodelling are found to be dominant. Details of the central values, the impact of\nindividual sources of uncertainties, and their correlations between experiments at\n\\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ can be found in Appendix~\\ref{app:tch}. \n\n\n\nThe shift in the combined cross-section at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ from a variation of $\\pm 1~\\GeV$ in the top-quark mass is $\\mp 0.8$~pb, which is similar to the shifts in the input measurements for the same $m_{\\myt}\\xspace$ variation. The shift in the combined cross-section at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ is not evaluated since no estimate is available for the CMS input measurement at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}.\n\n\\subsection{Combinations of {$\\boldmath \\ensuremath{tW}\\xspace$} cross-section measurements}\n\\label{sec:xscomb_tW}\nThe combination of the ATLAS and CMS $\\ensuremath{tW}\\xspace$ cross-section measurements\nat \\ensuremath{\\sqrt{s}=\\lhcenergySeven}~\\cite{TOPQ-2011-17, CMS-TOP-11-022} yields, after\ntwo iterations, a cross-section of \n\\begin{eqnarray*}\n\\ensuremath{\\sigmatW = 16.3\\pm2.3\\;(\\text{stat.})\\pm3.3\\;(\\text{syst.})\\pm0.7\\;(\\text{lumi.})\\,\\,\\text{pb} = 16.3\\pm4.1\\,\\,\\text{pb}}.\n\\end{eqnarray*}\nThe relative uncertainty is $\\ensuremath{25\\%}$, which improves on the uncertainty of 28\\% in the most precise individual\nmeasurement from CMS~\\cite{CMS-TOP-11-022}. The $\\chi^2$ for the combination is\n$\\ensuremath{0.01}$, corresponding to a probability of\n$\\ensuremath{91}\\%$. The CMS weight in the combination is\n$\\ensuremath{0.59}$, while the ATLAS weight is\n$\\ensuremath{0.41}$. The overall correlation between the\ntwo measurements is \\ensuremath{17}\\%. The contribution from\neach uncertainty category to the total uncertainty in the combined \\ensuremath{tW}\\xspace\\ \ncross-section measurement at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ is shown in\nTable~\\ref{tab:tWRes}(a).\n\nThe combination of the ATLAS and CMS $\\ensuremath{tW}\\xspace$ cross-section measurements\nat \\ensuremath{\\sqrt{s}=\\lhcenergyEight}~\\cite{TOPQ-2012-20,CMS-TOP-12-040} results, after\ntwo iterations, in\n\\begin{eqnarray*}\n\\ensuremath{\\sigmatW = 23.1\\pm1.1\\;(\\text{stat.})\\pm3.3\\;(\\text{syst.})\\pm0.8\\;(\\text{lumi.})\\,\\,\\text{pb} = 23.1\\pm3.6\\,\\,\\text{pb}}.\n\\end{eqnarray*}\nThe relative uncertainty is $\\ensuremath{15.6\\%}$, which improves on the uncertainty of 16.5\\% in the most precise individual\nmeasurement from ATLAS~\\cite{TOPQ-2012-20}. The $\\chi^2$ for the combination is\n$\\ensuremath{0.01}$, corresponding to a probability of\n$\\ensuremath{94}\\%$. The ATLAS weight in the combination is\n$\\ensuremath{0.70}$, while the CMS weight is\n$\\ensuremath{0.30}$. The overall correlation between the\ntwo measurements is \\ensuremath{40}\\%.\nSimilar to the \\myt-channel\\xspace, this is larger than the correlation between the\nmeasurements at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ due to the increased importance of the\ntheory modelling uncertainties.\nThe contribution from\neach uncertainty category to the total uncertainty in the combined \\ensuremath{tW}\\xspace\\ \ncross-section measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ is shown in\nTable~\\ref{tab:tWRes}(b).\n\n\\begin{table}[!htbp]\n \\caption{Contribution from each uncertainty category to the combined\n \\ensuremath{tW}\\xspace cross-section ($\\ensuremath{\\sigma_{\\mytW}}\\xspace$) uncertainty at\n (a) \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ and (b) \\ensuremath{\\sqrt{s}=\\lhcenergyEight}. The total\n uncertainty is computed by adding in quadrature all the individual systematic\n uncertainties (including the uncertainty in the integrated luminosity) and the statistical uncertainty in data.\n Correlations of systematic uncertainties between\n experiments are presented in Appendix~\\ref{app:tw}.}\n \\begin{center}\n \\begin{minipage}[b]{0.46\\hsize}\\centering\n \\captionof{subfigure}{}\n \\begin{tabular}{l | r | r}\n\\hline\n \\hline\n \\multicolumn{3}{c}{ $\\ensuremath{\\sigma_{\\mytW}}\\xspace$, $\\sqrt{s} = $ 7~\\TeV } \\\\\n\\hline\n\\hline\n{\\bf Combined cross-section} & \\multicolumn{2}{c}{16.3~pb} \\\\\n\\hline\n\\hline\n\\multirow{2}{*}{Uncertainty category} & \\multicolumn{2}{c}{Uncertainty} \\\\\\cline{2-3}\n & ~~[\\%] & [pb] \\\\\n\\hline\nData statistical & 14.0 & 2.3\\\\\n\\hline\nSimulation statistical & 0.8 & 0.1\\\\\nIntegrated luminosity & 4.4 & 0.7\\\\\nTheory modelling & 13.9 & 2.3\\\\\nBackground normalisation & 6.0 & 1.0\\\\\nJets & 11.5 & 1.9\\\\\nDetector modelling & 6.2 & 1.0\\\\\n\\hline\nTotal syst.\\ unc.\\ (excl.\\ lumi.) & 20.0 & 3.3\\\\\nTotal syst.\\ unc.\\ (incl.\\ lumi.) & 20.5 & 3.3\\\\\n\\hline\\hline\n{\\bf Total uncertainty} & 24.8 & 4.1\\\\\n\\hline\n \\hline\n\\end{tabular}\n \\end{minipage}\n \\hfill\n \\begin{minipage}[b]{0.46\\hsize}\\centering\n \\captionof{subfigure}{}\n \\begin{tabular}{l | r | r}\n\\hline\n \\hline\n \\multicolumn{3}{c}{ $\\ensuremath{\\sigma_{\\mytW}}\\xspace$, $\\sqrt{s} = $ 8~\\TeV } \\\\\n\\hline\n\\hline\n{\\bf Combined cross-section} & \\multicolumn{2}{c}{23.1~pb} \\\\\n\\hline\n\\hline\n\\multirow{2}{*}{Uncertainty category} & \\multicolumn{2}{c}{Uncertainty} \\\\\\cline{2-3}\n & ~~[\\%] & [pb] \\\\\n\\hline\nData statistical & 4.7 & 1.1\\\\\n\\hline\nSimulation statistical & 0.8 & 0.2\\\\\nIntegrated luminosity & 3.6 & 0.8\\\\\nTheory modelling & 11.8 & 2.7\\\\\nBackground normalisation & 2.2 & 0.5\\\\\nJets & 6.2 & 1.4\\\\\nDetector modelling & 4.9 & 1.1\\\\\n\\hline\nTotal syst.\\ unc.\\ (excl.\\ lumi.) & 14.4 & 3.3\\\\\nTotal syst.\\ unc.\\ (incl.\\ lumi.) & 14.8 & 3.4\\\\\n\\hline\\hline\n{\\bf Total uncertainty} & 15.6 & 3.6\\\\\n\\hline\n \\hline\n\\end{tabular}\n \\end{minipage}\n \\end{center}\n \\label{tab:tWRes}\n\\end{table}\n\nAt both centre-of-mass energies, the uncertainties in the theory\nmodelling are found to be dominant. The jet uncertainties are also\nimportant. Details of the central values, the impact of\nindividual sources of uncertainties, and their correlations between experiments at\n\\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ are presented in Appendix~\\ref{app:tw}.\n\nThe shift in the combined cross-section at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ from a variation of $\\pm\n1~\\GeV$ in the top-quark mass is $\\pm 1.1$~pb, which is similar in magnitude to that in the input measurements for the same $m_{\\myt}\\xspace$ variation. The shift in the combined cross-section at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}\\ is not evaluated since no estimates are available for the input measurements at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}.\n\n\\subsection{Combination of \\mys-channel\\xspace\\ cross-section measurements}\n\\label{sec:xscomb_sch}\n\nThe ATLAS and CMS \\mys-channel\\xspace\\ cross-section measurements suffer from large backgrounds, and the cross-section measurements have large uncertainties. Since the systematic uncertainties mainly affect the background prediction, they are not scaled in the iterative BLUE procedure. Only the luminosity uncertainty is scaled with the central value. The combination of the ATLAS and CMS \\mys-channel\\xspace\\ cross-section measurements at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}~\\cite{TOPQ-2015-01, CMS-TOP-13-009} results, after\ntwo iterations, in a cross-section of\n\n\\begin{eqnarray*}\n\\ensuremath{\\sigmasch = 4.9\\pm0.8\\;(\\text{stat.})\\pm1.2\\;(\\text{syst.})\\pm0.2\\;(\\text{lumi.})\\,\\,\\text{pb} = 4.9\\pm1.4\\,\\,\\text{pb}}.\n\\end{eqnarray*}\nThe relative uncertainty is $\\ensuremath{30\\%}$, very similar to the most precise individual measurement from ATLAS~\\cite{TOPQ-2015-01}. The $\\chi^2$ for the combination is $\\ensuremath{1.45}$, corresponding to a probability of $\\ensuremath{23}\\%$. \nThe ATLAS weight in the combination is $\\ensuremath{0.99}$, while the CMS weight is $\\ensuremath{0.01}$. The overall correlation between the two measurements is \\ensuremath{15}\\%. The contribution from each uncertainty category to the total uncertainty in the combined \\mys-channel\\xspace\\ cross-section measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ is shown in Table~\\ref{tab:sChanRes}.\n\n\\begin{table}[!htbp]\n \\caption{Contribution from each uncertainty category to the combined\n \\mys-channel\\xspace\\ cross-section ($\\ensuremath{\\sigma_{\\mys\\textrm{-chan.}}}\\xspace$) uncertainty at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}. The total\n uncertainty is computed by adding in quadrature all the individual systematic\n uncertainties (including the uncertainty in the integrated luminosity) and the statistical uncertainty in data.\n Correlations of systematic uncertainties between\n experiments are presented in Appendix~\\ref{app:sch}.}\n \\begin{center}\n \\begin{tabular}{l | r | r}\n\\hline\n \\hline\n \\multicolumn{3}{c}{ $\\ensuremath{\\sigma_{\\mys\\textrm{-chan.}}}\\xspace$, $\\sqrt{s} = $ 8~\\TeV } \\\\\n\\hline\n\\hline\n{\\bf Combined cross-section} & \\multicolumn{2}{c}{4.9~pb} \\\\\n\\hline\n\\hline\n\\multirow{2}{*}{Uncertainty category} & \\multicolumn{2}{c}{Uncertainty} \\\\\\cline{2-3}\n & ~~[\\%] & [pb] \\\\\n\\hline\nData statistical & 16 & 0.8\\\\\n\\hline\nSimulation statistical & 12 & 0.6\\\\\nIntegrated luminosity & 5 & 0.2\\\\\nTheory modelling & 14 & 0.7\\\\\nBackground normalisation & 8 & 0.4\\\\\nJets & 13 & 0.6\\\\\nDetector modelling & 8 & 0.4\\\\\n\\hline\nTotal syst.\\ unc.\\ (excl.\\ lumi.) & 25 & 1.2\\\\\nTotal syst.\\ unc.\\ (incl.\\ lumi.) & 25 & 1.2\\\\\n\\hline\\hline\n{\\bf Total uncertainty} & 30 & 1.4\\\\\n\\hline\n \\hline\n\\end{tabular}\n \\end{center}\n \\label{tab:sChanRes}\n\\end{table}\n\nSince the ATLAS measurement has a large weight in the combination, the importance of each uncertainty in the combination is similar to that in the ATLAS measurement, as presented in Appendix~\\ref{app:sch}.\n\nThe shift in the combined cross-section at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ from a variation in the top-quark mass is not evaluated since no estimate is available for the ATLAS input measurement.\n\n\\subsection{Summary of cross-section combinations}\nA summary of the cross-sections measured by ATLAS and CMS\nand their combinations in all single-top-quark production modes at each\ncentre-of-mass energy is shown in\nFigure~\\ref{fig:xssum}. The measurements are compared with the theoretical predictions shown in Table~\\ref{tab:xs}: NNLO for \\myt-channel\\xspace\\ only, NLO and NLO+NNLL for all three production modes. For the NLO calculation, the renormalisation- and factorisation-scale uncertainties and the sum in quadrature of the contributions from scale, PDF, and $\\alpha_{\\rm s}$ are shown separately. Only the scale uncertainty is shown for the NNLO calculation. \nFor the NLO+NNLL calculation, the sum in quadrature of the contributions from scale, PDF, and $\\alpha_{\\rm s}$ is shown. All measurements are in good agreement with their corresponding theoretical predictions within their total uncertainties.\n\n\n\\begin{figure}[!ht]\n \\begin{center}\n \\includegraphics[width=1.0\\textwidth]{xs\/singletop_xsec_combinations}\n \\end{center}\n \\caption{Single-top-quark cross-section measurements performed by ATLAS and CMS, together with the combined results shown in Sections~\\ref{sec:xscomb_tch}$-$\\ref{sec:xscomb_sch}. These measurements are compared with the theoretical predictions at NLO and NLO+NNLL for all three production modes and the prediction at NNLO for \\myt-channel\\xspace\\ only. The corresponding theoretical uncertainties are also presented. The scale uncertainty for the NNLO prediction is small and is presented as a narrow band under the dashed line.}\n \\label{fig:xssum}\n\\end{figure}\n\nThe stability of the combinations of the cross-section measurements to variations in\nthe correlation assumptions, discussed in Section~\\ref{sec:systcat}, is checked for the theory modelling, JES, the most important contributions to the theoretical cross-section predictions\n(i.e.\\ PDF+$\\alpha_{\\rm s}$ and scale) and the integrated luminosity. The results\nof these tests show that their impacts on the cross-section\ncombinations are very small, similar to the stability tests\nfor the combination of the $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ values discussed in Section~\\ref{sec:tests}.\n\n\n\n\n\n \\section{Combinations of {\\boldmath $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$} determinations}\n\\label{sec:vtbcomb}\n\\newcommand{\\ensuremath{0.03}}{\\ensuremath{0.03}}\n\\newcommand{\\ensuremath{1.02}}{\\ensuremath{1.02}}\n\\newcommand{\\ensuremath{1.04}}{\\ensuremath{1.04}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 1.04\\pm0.01\\;(\\text{stat.})\\pm0.06\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.}) = 1.04\\pm0.07}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 1.04\\pm0.01\\;(\\text{stat.})\\pm0.06\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.}) = 1.04\\pm0.07}}\n\\newcommand{\\ensuremath{6.5\\%}}{\\ensuremath{6.5\\%}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 1.02\\pm0.01\\;(\\text{stat.})\\pm0.03\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.}) = 1.02\\pm0.03}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 1.02\\pm0.01\\;(\\text{stat.})\\pm0.03\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.}) = 1.02\\pm0.03}}\n\\newcommand{\\ensuremath{1.02\\pm0.03}}{\\ensuremath{1.02\\pm0.03}}\n\\newcommand{\\ensuremath{3.2\\%}}{\\ensuremath{3.2\\%}}\n\\newcommand{\\ensuremath{0.63}}{\\ensuremath{0.63}}\n\\newcommand{\\ensuremath{89}}{\\ensuremath{89}}\n\\newcommand{\\ensuremath{0.09}}{\\ensuremath{0.09}}\n\\newcommand{\\ensuremath{0.25}}{\\ensuremath{0.25}}\n\\newcommand{\\ensuremath{0.48}}{\\ensuremath{0.48}}\n\\newcommand{\\ensuremath{0.53}}{\\ensuremath{0.53}}\n\\newcommand{\\ensuremath{0.16}}{\\ensuremath{0.16}}\n\\newcommand{\\ensuremath{0.16}}{\\ensuremath{0.16}}\n\\newcommand{\\ensuremath{0.27}}{\\ensuremath{0.27}}\n\\newcommand{\\ensuremath{-0.80}}{\\ensuremath{-0.80}}\n\\newcommand{\\ensuremath{0.42}}{\\ensuremath{0.42}}\n\\newcommand{\\ensuremath{0.46}}{\\ensuremath{0.46}}\n\\newcommand{\\ensuremath{0.40}}{\\ensuremath{0.40}}\n\\newcommand{\\ensuremath{0.21}}{\\ensuremath{0.21}}\n\\newcommand{\\ensuremath{0.33}}{\\ensuremath{0.33}}\n\\newcommand{\\ensuremath{0.20}}{\\ensuremath{0.20}}\n\\newcommand{\\ensuremath{1.56}}{\\ensuremath{1.56}}\n \\newcommand{\\ensuremath{0.08}}{\\ensuremath{0.08}}\n\\newcommand{\\ensuremath{1.02}}{\\ensuremath{1.02}}\n\\newcommand{\\ensuremath{1.03}}{\\ensuremath{1.03}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 1.03\\pm0.05\\;(\\text{stat.})\\pm0.14\\;(\\text{syst.})\\pm0.03\\;(\\text{lumi.}) = 1.03\\pm0.16}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 1.03\\pm0.05\\;(\\text{stat.})\\pm0.14\\;(\\text{syst.})\\pm0.03\\;(\\text{lumi.}) = 1.03\\pm0.16}}\n\\newcommand{\\ensuremath{15.2\\%}}{\\ensuremath{15.2\\%}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 1.02\\pm0.03\\;(\\text{stat.})\\pm0.07\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.}) = 1.02\\pm0.08}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 1.02\\pm0.03\\;(\\text{stat.})\\pm0.07\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.}) = 1.02\\pm0.08}}\n\\newcommand{\\ensuremath{1.02\\pm0.08}}{\\ensuremath{1.02\\pm0.08}}\n\\newcommand{\\ensuremath{7.6\\%}}{\\ensuremath{7.6\\%}}\n\\newcommand{\\ensuremath{0.02}}{\\ensuremath{0.02}}\n\\newcommand{\\ensuremath{100}}{\\ensuremath{100}}\n\\newcommand{\\ensuremath{0.08}}{\\ensuremath{0.08}}\n\\newcommand{\\ensuremath{0.11}}{\\ensuremath{0.11}}\n\\newcommand{\\ensuremath{0.61}}{\\ensuremath{0.61}}\n\\newcommand{\\ensuremath{-0.05}}{\\ensuremath{-0.05}}\n\\newcommand{\\ensuremath{0.16}}{\\ensuremath{0.16}}\n\\newcommand{\\ensuremath{-0.06}}{\\ensuremath{-0.06}}\n\\newcommand{\\ensuremath{0.15}}{\\ensuremath{0.15}}\n\\newcommand{\\ensuremath{0.08}}{\\ensuremath{0.08}}\n\\newcommand{\\ensuremath{0.43}}{\\ensuremath{0.43}}\n\\newcommand{\\ensuremath{0.15}}{\\ensuremath{0.15}}\n\\newcommand{\\ensuremath{0.12}}{\\ensuremath{0.12}}\n\\newcommand{\\ensuremath{0.35}}{\\ensuremath{0.35}}\n\\newcommand{\\ensuremath{0.29}}{\\ensuremath{0.29}}\n\\newcommand{\\ensuremath{0.17}}{\\ensuremath{0.17}}\n\\newcommand{\\ensuremath{1.55}}{\\ensuremath{1.55}}\n \\newcommand{\\ensuremath{0.14}}{\\ensuremath{0.14}}\n\\newcommand{\\ensuremath{0.97}}{\\ensuremath{0.97}}\n\\newcommand{\\ensuremath{0.93}}{\\ensuremath{0.93}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 0.93\\pm0.15\\;(\\text{stat.})\\pm0.23\\;(\\text{syst.})\\pm0.05\\;(\\text{lumi.}) = 0.93\\pm0.28}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 0.93\\pm0.15\\;(\\text{stat.})\\pm0.23\\;(\\text{syst.})\\pm0.05\\;(\\text{lumi.}) = 0.93\\pm0.28}}\n\\newcommand{\\ensuremath{29.7\\%}}{\\ensuremath{29.7\\%}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 0.97\\pm0.08\\;(\\text{stat.})\\pm0.12\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.}) = 0.97\\pm0.14}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 0.97\\pm0.08\\;(\\text{stat.})\\pm0.12\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.}) = 0.97\\pm0.14}}\n\\newcommand{\\ensuremath{0.97\\pm0.14}}{\\ensuremath{0.97\\pm0.14}}\n\\newcommand{\\ensuremath{14.8\\%}}{\\ensuremath{14.8\\%}}\n\\newcommand{\\ensuremath{1.42}}{\\ensuremath{1.42}}\n\\newcommand{\\ensuremath{23}}{\\ensuremath{23}}\n\\newcommand{\\ensuremath{0.99}}{\\ensuremath{0.99}}\n\\newcommand{\\ensuremath{-1.19}}{\\ensuremath{-1.19}}\n\\newcommand{\\ensuremath{0.01}}{\\ensuremath{0.01}}\n\\newcommand{\\ensuremath{1.19}}{\\ensuremath{1.19}}\n\\newcommand{\\ensuremath{0.15}}{\\ensuremath{0.15}}\n\\newcommand{\\ensuremath{3.60}}{\\ensuremath{3.60}}\n \\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 1.04\\pm0.03\\;(\\text{stat.})\\pm0.04\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 1.04\\pm0.05}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 1.04\\pm0.03\\;(\\text{stat.})\\pm0.04\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 1.04\\pm0.05}}\n\\newcommand{\\ensuremath{1.04}}{\\ensuremath{1.04}}\n\\newcommand{\\ensuremath{0.03}}{\\ensuremath{0.03}}\n\\newcommand{\\ensuremath{0.04}}{\\ensuremath{0.04}}\n\\newcommand{\\ensuremath{0.05}}{\\ensuremath{0.05}}\n\\newcommand{\\ensuremath{0.02}}{\\ensuremath{0.02}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 1.02\\pm0.01\\;(\\text{stat.})\\pm0.02\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 1.02\\pm0.02}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 1.02\\pm0.01\\;(\\text{stat.})\\pm0.02\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 1.02\\pm0.02}}\n\\newcommand{\\ensuremath{1.02\\pm0.02}}{\\ensuremath{1.02\\pm0.02}}\n\\newcommand{\\ensuremath{2.4\\%}}{\\ensuremath{2.4\\%}}\n\\newcommand{\\ensuremath{1.54}}{\\ensuremath{1.54}}\n\\newcommand{\\ensuremath{67}}{\\ensuremath{67}}\n\\newcommand{\\ensuremath{0.23}}{\\ensuremath{0.23}}\n\\newcommand{\\ensuremath{0.58}}{\\ensuremath{0.58}}\n\\newcommand{\\ensuremath{0.26}}{\\ensuremath{0.26}}\n\\newcommand{\\ensuremath{0.46}}{\\ensuremath{0.46}}\n\\newcommand{\\ensuremath{0.23}}{\\ensuremath{0.23}}\n\\newcommand{\\ensuremath{0.31}}{\\ensuremath{0.31}}\n\\newcommand{\\ensuremath{0.28}}{\\ensuremath{0.28}}\n\\newcommand{\\ensuremath{-1.23}}{\\ensuremath{-1.23}}\n\\newcommand{\\ensuremath{0.43}}{\\ensuremath{0.43}}\n\\newcommand{\\ensuremath{0.44}}{\\ensuremath{0.44}}\n\\newcommand{\\ensuremath{0.44}}{\\ensuremath{0.44}}\n\\newcommand{\\ensuremath{0.44}}{\\ensuremath{0.44}}\n\\newcommand{\\ensuremath{0.44}}{\\ensuremath{0.44}}\n\\newcommand{\\ensuremath{0.45}}{\\ensuremath{0.45}}\n\\newcommand{\\ensuremath{1.78}}{\\ensuremath{1.78}}\n \\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 1.04\\pm0.03\\;(\\text{stat.})\\pm0.07\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 1.04\\pm0.08}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 1.04\\pm0.03\\;(\\text{stat.})\\pm0.07\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 1.04\\pm0.08}}\n\\newcommand{\\ensuremath{1.04}}{\\ensuremath{1.04}}\n\\newcommand{\\ensuremath{0.03}}{\\ensuremath{0.03}}\n\\newcommand{\\ensuremath{0.07}}{\\ensuremath{0.07}}\n\\newcommand{\\ensuremath{0.08}}{\\ensuremath{0.08}}\n\\newcommand{\\ensuremath{0.04}}{\\ensuremath{0.04}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 1.02\\pm0.01\\;(\\text{stat.})\\pm0.04\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 1.02\\pm0.04}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 1.02\\pm0.01\\;(\\text{stat.})\\pm0.04\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 1.02\\pm0.04}}\n\\newcommand{\\ensuremath{1.02\\pm0.04}}{\\ensuremath{1.02\\pm0.04}}\n\\newcommand{\\ensuremath{3.8\\%}}{\\ensuremath{3.8\\%}}\n\\newcommand{\\ensuremath{0.57}}{\\ensuremath{0.57}}\n\\newcommand{\\ensuremath{90}}{\\ensuremath{90}}\n\\newcommand{\\ensuremath{0.08}}{\\ensuremath{0.08}}\n\\newcommand{\\ensuremath{0.58}}{\\ensuremath{0.58}}\n\\newcommand{\\ensuremath{0.38}}{\\ensuremath{0.38}}\n\\newcommand{\\ensuremath{-0.20}}{\\ensuremath{-0.20}}\n\\newcommand{\\ensuremath{0.18}}{\\ensuremath{0.18}}\n\\newcommand{\\ensuremath{-0.38}}{\\ensuremath{-0.38}}\n\\newcommand{\\ensuremath{0.36}}{\\ensuremath{0.36}}\n\\newcommand{\\ensuremath{0.27}}{\\ensuremath{0.27}}\n\\newcommand{\\ensuremath{0.68}}{\\ensuremath{0.68}}\n\\newcommand{\\ensuremath{0.70}}{\\ensuremath{0.70}}\n\\newcommand{\\ensuremath{0.69}}{\\ensuremath{0.69}}\n\\newcommand{\\ensuremath{0.70}}{\\ensuremath{0.70}}\n\\newcommand{\\ensuremath{0.70}}{\\ensuremath{0.70}}\n\\newcommand{\\ensuremath{0.72}}{\\ensuremath{0.72}}\n\\newcommand{\\ensuremath{1.70}}{\\ensuremath{1.70}}\n \\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 0.97\\pm0.05\\;(\\text{stat.})\\pm0.04\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 0.97\\pm0.07}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 0.97\\pm0.05\\;(\\text{stat.})\\pm0.04\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 0.97\\pm0.07}}\n\\newcommand{\\ensuremath{0.97}}{\\ensuremath{0.97}}\n\\newcommand{\\ensuremath{0.05}}{\\ensuremath{0.05}}\n\\newcommand{\\ensuremath{0.04}}{\\ensuremath{0.04}}\n\\newcommand{\\ensuremath{0.07}}{\\ensuremath{0.07}}\n\\newcommand{\\ensuremath{0.03}}{\\ensuremath{0.03}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 0.98\\pm0.02\\;(\\text{stat.})\\pm0.02\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 0.98\\pm0.03}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 0.98\\pm0.02\\;(\\text{stat.})\\pm0.02\\;(\\text{syst.})\\pm0.00\\;(\\text{lumi.}) = 0.98\\pm0.03}}\n\\newcommand{\\ensuremath{0.98\\pm0.03}}{\\ensuremath{0.98\\pm0.03}}\n\\newcommand{\\ensuremath{3.4\\%}}{\\ensuremath{3.4\\%}}\n\\newcommand{\\ensuremath{300.77}}{\\ensuremath{300.77}}\n\\newcommand{\\ensuremath{0}}{\\ensuremath{0}}\n\\newcommand{\\ensuremath{0.97}}{\\ensuremath{0.97}}\n\\newcommand{\\ensuremath{-17.34}}{\\ensuremath{-17.34}}\n\\newcommand{\\ensuremath{0.03}}{\\ensuremath{0.03}}\n\\newcommand{\\ensuremath{17.34}}{\\ensuremath{17.34}}\n\\newcommand{\\ensuremath{0.55}}{\\ensuremath{0.55}}\n\\newcommand{\\ensuremath{18.56}}{\\ensuremath{18.56}}\n \\newcommand{\\ensuremath{0.04}}{\\ensuremath{0.04}}\n\\newcommand{\\ensuremath{0.02}}{\\ensuremath{0.02}}\n\\newcommand{\\ensuremath{0.04}}{\\ensuremath{0.04}}\n\\newcommand{\\ensuremath{1.02}}{\\ensuremath{1.02}}\n\\newcommand{\\ensuremath{1.04}}{\\ensuremath{1.04}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 1.04\\pm0.01\\;(\\text{stat.})\\pm0.08\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.}) = 1.04\\pm0.08}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 1.04\\pm0.01\\;(\\text{stat.})\\pm0.08\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.}) = 1.04\\pm0.08}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 1.04\\pm0.01\\;(\\text{stat.})\\pm0.06\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.})\\pm0.04\\;(\\text{theo.}) = 1.04\\pm0.08}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 1.04\\pm0.01\\;(\\text{stat.})\\pm0.06\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.})\\pm0.04\\;(\\text{theo.}) = 1.04\\pm0.08}}\n\\newcommand{\\ensuremath{4.2\\%}}{\\ensuremath{4.2\\%}}\n\\newcommand{\\ensuremath{2.1\\%}}{\\ensuremath{2.1\\%}}\n\\newcommand{\\ensuremath{7.7\\%}}{\\ensuremath{7.7\\%}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 1.02\\pm0.01\\;(\\text{stat.})\\pm0.04\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.}) = 1.02\\pm0.04}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 1.02\\pm0.01\\;(\\text{stat.})\\pm0.04\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.}) = 1.02\\pm0.04}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|=~&~1.02\\pm0.01\\;(\\text{stat.})\\pm0.03\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.})\\pm0.02\\;(\\text{theo.}) \\\\=~&~1.02\\pm0.04\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.}) = 1.02\\pm0.04}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|=~&~1.02\\pm0.01\\;(\\text{stat.})\\pm0.03\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.})\\pm0.02\\;(\\text{theo.}) \\\\=~&~1.02\\pm0.04\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.}) = 1.02\\pm0.04}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 1.02\\pm0.04\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.})}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 1.02\\pm0.04\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.})}}\n\\newcommand{\\ensuremath{1.02\\pm0.04}}{\\ensuremath{1.02\\pm0.04}}\n\\newcommand{\\ensuremath{3.9\\%}}{\\ensuremath{3.9\\%}}\n\\newcommand{\\ensuremath{0.63}}{\\ensuremath{0.63}}\n\\newcommand{\\ensuremath{89}}{\\ensuremath{89}}\n\\newcommand{\\ensuremath{0.09}}{\\ensuremath{0.09}}\n\\newcommand{\\ensuremath{0.25}}{\\ensuremath{0.25}}\n\\newcommand{\\ensuremath{0.49}}{\\ensuremath{0.49}}\n\\newcommand{\\ensuremath{0.54}}{\\ensuremath{0.54}}\n\\newcommand{\\ensuremath{0.15}}{\\ensuremath{0.15}}\n\\newcommand{\\ensuremath{0.16}}{\\ensuremath{0.16}}\n\\newcommand{\\ensuremath{0.27}}{\\ensuremath{0.27}}\n\\newcommand{\\ensuremath{-0.80}}{\\ensuremath{-0.80}}\n\\newcommand{\\ensuremath{0.54}}{\\ensuremath{0.54}}\n\\newcommand{\\ensuremath{0.55}}{\\ensuremath{0.55}}\n\\newcommand{\\ensuremath{0.52}}{\\ensuremath{0.52}}\n\\newcommand{\\ensuremath{0.33}}{\\ensuremath{0.33}}\n\\newcommand{\\ensuremath{0.44}}{\\ensuremath{0.44}}\n\\newcommand{\\ensuremath{0.31}}{\\ensuremath{0.31}}\n\\newcommand{\\ensuremath{1.45}}{\\ensuremath{1.45}}\n \\newcommand{\\ensuremath{0.08}}{\\ensuremath{0.08}}\n\\newcommand{\\ensuremath{0.04}}{\\ensuremath{0.04}}\n\\newcommand{\\ensuremath{0.09}}{\\ensuremath{0.09}}\n\\newcommand{\\ensuremath{1.02}}{\\ensuremath{1.02}}\n\\newcommand{\\ensuremath{1.03}}{\\ensuremath{1.03}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 1.03\\pm0.05\\;(\\text{stat.})\\pm0.16\\;(\\text{syst.})\\pm0.03\\;(\\text{lumi.}) = 1.03\\pm0.17}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 1.03\\pm0.05\\;(\\text{stat.})\\pm0.16\\;(\\text{syst.})\\pm0.03\\;(\\text{lumi.}) = 1.03\\pm0.17}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 1.03\\pm0.05\\;(\\text{stat.})\\pm0.14\\;(\\text{syst.})\\pm0.03\\;(\\text{lumi.})\\pm0.07\\;(\\text{theo.}) = 1.03\\pm0.17}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 1.03\\pm0.05\\;(\\text{stat.})\\pm0.14\\;(\\text{syst.})\\pm0.03\\;(\\text{lumi.})\\pm0.07\\;(\\text{theo.}) = 1.03\\pm0.17}}\n\\newcommand{\\ensuremath{7.2\\%}}{\\ensuremath{7.2\\%}}\n\\newcommand{\\ensuremath{3.6\\%}}{\\ensuremath{3.6\\%}}\n\\newcommand{\\ensuremath{16.8\\%}}{\\ensuremath{16.8\\%}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 1.02\\pm0.03\\;(\\text{stat.})\\pm0.08\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.}) = 1.02\\pm0.09}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 1.02\\pm0.03\\;(\\text{stat.})\\pm0.08\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.}) = 1.02\\pm0.09}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|=~&~1.02\\pm0.03\\;(\\text{stat.})\\pm0.07\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.})\\pm0.04\\;(\\text{theo.}) \\\\=~&~1.02\\pm0.09\\;(\\text{meas.})\\pm0.04\\;(\\text{theo.}) = 1.02\\pm0.09}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|=~&~1.02\\pm0.03\\;(\\text{stat.})\\pm0.07\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.})\\pm0.04\\;(\\text{theo.}) \\\\=~&~1.02\\pm0.09\\;(\\text{meas.})\\pm0.04\\;(\\text{theo.}) = 1.02\\pm0.09}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 1.02\\pm0.09\\;(\\text{meas.})\\pm0.04\\;(\\text{theo.})}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 1.02\\pm0.09\\;(\\text{meas.})\\pm0.04\\;(\\text{theo.})}}\n\\newcommand{\\ensuremath{1.02\\pm0.09}}{\\ensuremath{1.02\\pm0.09}}\n\\newcommand{\\ensuremath{8.4\\%}}{\\ensuremath{8.4\\%}}\n\\newcommand{\\ensuremath{0.02}}{\\ensuremath{0.02}}\n\\newcommand{\\ensuremath{100}}{\\ensuremath{100}}\n\\newcommand{\\ensuremath{0.08}}{\\ensuremath{0.08}}\n\\newcommand{\\ensuremath{0.11}}{\\ensuremath{0.11}}\n\\newcommand{\\ensuremath{0.61}}{\\ensuremath{0.61}}\n\\newcommand{\\ensuremath{-0.05}}{\\ensuremath{-0.05}}\n\\newcommand{\\ensuremath{0.15}}{\\ensuremath{0.15}}\n\\newcommand{\\ensuremath{-0.05}}{\\ensuremath{-0.05}}\n\\newcommand{\\ensuremath{0.16}}{\\ensuremath{0.16}}\n\\newcommand{\\ensuremath{0.08}}{\\ensuremath{0.08}}\n\\newcommand{\\ensuremath{0.49}}{\\ensuremath{0.49}}\n\\newcommand{\\ensuremath{0.22}}{\\ensuremath{0.22}}\n\\newcommand{\\ensuremath{0.20}}{\\ensuremath{0.20}}\n\\newcommand{\\ensuremath{0.39}}{\\ensuremath{0.39}}\n\\newcommand{\\ensuremath{0.34}}{\\ensuremath{0.34}}\n\\newcommand{\\ensuremath{0.21}}{\\ensuremath{0.21}}\n\\newcommand{\\ensuremath{1.54}}{\\ensuremath{1.54}}\n \\newcommand{\\ensuremath{0.12}}{\\ensuremath{0.12}}\n\\newcommand{\\ensuremath{0.02}}{\\ensuremath{0.02}}\n\\newcommand{\\ensuremath{0.15}}{\\ensuremath{0.15}}\n\\newcommand{\\ensuremath{0.97}}{\\ensuremath{0.97}}\n\\newcommand{\\ensuremath{0.93}}{\\ensuremath{0.93}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 0.93\\pm0.15\\;(\\text{stat.})\\pm0.23\\;(\\text{syst.})\\pm0.05\\;(\\text{lumi.}) = 0.93\\pm0.28}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 0.93\\pm0.15\\;(\\text{stat.})\\pm0.23\\;(\\text{syst.})\\pm0.05\\;(\\text{lumi.}) = 0.93\\pm0.28}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|^2 = 0.93\\pm0.15\\;(\\text{stat.})\\pm0.23\\;(\\text{syst.})\\pm0.05\\;(\\text{lumi.})\\pm0.04\\;(\\text{theo.}) = 0.93\\pm0.28}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|^2 = 0.93\\pm0.15\\;(\\text{stat.})\\pm0.23\\;(\\text{syst.})\\pm0.05\\;(\\text{lumi.})\\pm0.04\\;(\\text{theo.}) = 0.93\\pm0.28}}\n\\newcommand{\\ensuremath{4.6\\%}}{\\ensuremath{4.6\\%}}\n\\newcommand{\\ensuremath{2.3\\%}}{\\ensuremath{2.3\\%}}\n\\newcommand{\\ensuremath{30.1\\%}}{\\ensuremath{30.1\\%}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 0.97\\pm0.08\\;(\\text{stat.})\\pm0.12\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.}) = 0.97\\pm0.15}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 0.97\\pm0.08\\;(\\text{stat.})\\pm0.12\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.}) = 0.97\\pm0.15}}\n\\newcommand{\\ensuremath{|\\flv\\vtb|=~&~0.97\\pm0.08\\;(\\text{stat.})\\pm0.12\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.})\\pm0.02\\;(\\text{theo.}) \\\\=~&~0.97\\pm0.15\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.}) = 0.97\\pm0.15}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace|=~&~0.97\\pm0.08\\;(\\text{stat.})\\pm0.12\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.})\\pm0.02\\;(\\text{theo.}) \\\\=~&~0.97\\pm0.15\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.}) = 0.97\\pm0.15}}\n\\newcommand{\\ensuremath{|\\flv\\vtb| = 0.97\\pm0.15\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.})}}{\\ensuremath{|\\ensuremath{f_{\\rm LV}}\\xspace\\ensuremath{V_{\\myt\\myb}}\\xspace| = 0.97\\pm0.15\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.})}}\n\\newcommand{\\ensuremath{0.97\\pm0.15}}{\\ensuremath{0.97\\pm0.15}}\n\\newcommand{\\ensuremath{15.0\\%}}{\\ensuremath{15.0\\%}}\n\\newcommand{\\ensuremath{1.42}}{\\ensuremath{1.42}}\n\\newcommand{\\ensuremath{23}}{\\ensuremath{23}}\n\\newcommand{\\ensuremath{0.99}}{\\ensuremath{0.99}}\n\\newcommand{\\ensuremath{-1.19}}{\\ensuremath{-1.19}}\n\\newcommand{\\ensuremath{0.01}}{\\ensuremath{0.01}}\n\\newcommand{\\ensuremath{1.19}}{\\ensuremath{1.19}}\n\\newcommand{\\ensuremath{0.15}}{\\ensuremath{0.15}}\n\\newcommand{\\ensuremath{3.88}}{\\ensuremath{3.88}}\n The measured cross-section for a given single-top-quark production mode, $\\sigma_{\\rm meas.}$, has\na linear dependence on $\\ensuremath{\\absvtb^2}\\xspace$ as defined in Eq.~(\\ref{eq:vtb}). Thus, a value of $\\ensuremath{\\absvtb^2}\\xspace$ is extracted from each cross-section measurement and the corresponding theoretical prediction (presented in\nSections~\\ref{sec:atlascmsxs} and \\ref{sec:theory}\nrespectively). These values are then combined per channel, and in an\noverall $\\ensuremath{\\absvtb^2}\\xspace$ combination. In the overall combination, the value from the CMS measurement of $\\ensuremath{\\sigma_{\\mys\\textrm{-chan.}}}\\xspace$ is excluded. The reason for excluding the CMS \\mys-channel\\xspace\\ analysis from the overall $\\ensuremath{\\absvtb^2}\\xspace$ combination is that, at the same centre-of-mass energy, the CMS \\myt-channel\\xspace\\ determination has strong correlations with the \\mys-channel\\xspace\\ determination, which contains relatively large uncertainties. The strong\ncorrelation between these two measurements makes the combined $\\ensuremath{\\absvtb^2}\\xspace$ value strongly\ndependent on the correlation assumptions for the dominant uncertainties.\nThis results in a large variation of the combined $\\ensuremath{\\absvtb^2}\\xspace$ value for different correlation assumptions.\n\n\n\nAll uncertainties in $\\ensuremath{\\sigma_{\\rm meas.}}\\xspace$ and $\\ensuremath{\\sigma_{\\rm theo.}}\\xspace$ are propagated to the $\\ensuremath{\\absvtb^2}\\xspace$ values,\ntaking into account the correlations described\nin Section~\\ref{sec:systcat}. The combined value of $\\ensuremath{\\absvtb^2}\\xspace$ is evaluated using the reference theoretical\ncross-section central values marked with a $^\\dagger$ in\nTable~\\ref{tab:xs}, where it can also be seen that the $E_{\\rm\n beam}$ uncertainty is negligible compared to other\nuncertainties. For the most precise measurements (i.e.\\ for\n$\\ensuremath{\\sigma_{{\\myt\\textrm{-chan.}}}}\\xspace$ cross-section measurements at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}), which have a large expected impact on the combination,\nthe other theoretical calculations from Table~\\ref{tab:xs} are used as\ncross-checks.\n\nTable~\\ref{tab:uncVtb} contains a summary of the individual $\\ensuremath{\\absvtb^2}\\xspace$\ndeterminations that are the inputs to the overall $\\ensuremath{\\absvtb^2}\\xspace$ combination, together with their experimental and theoretical uncertainties using the reference theoretical cross-sections and uncertainties. For the same processes and at the same centre-of-mass\nenergies, there are some important differences between uncertainty\ncategories. In analyses based on \\myt-channel\\xspace\\ events at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, the data\nstatistical uncertainty is larger in CMS than in ATLAS because the two experiments use data samples of different size. Differences in the category of jet uncertainties are due to the evaluation of the JES uncertainty in ATLAS using pseudoexperiments, while this uncertainty is introduced as a nuisance parameter in the fit in CMS. At \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, the difference between ATLAS and CMS in the background-normalisation category is due to the different techniques used to estimate each background\nuncertainty. Additional details are discussed in Appendix~\\ref{app:tch}. In the CMS\n\\ensuremath{tW}\\xspace\\ analysis at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, the uncertainty\nassociated with the size of the simulated samples is evaluated as part\nof the total statistical uncertainty. The large difference in the\npile-up uncertainty between ATLAS and CMS is due to the different\nmethods used to assess this uncertainty, as discussed in\nSection~\\ref{sec:systcat_exp}.\nAt \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, the sizes of the data and simulated samples used in the CMS \\ensuremath{tW}\\xspace\\ analysis are smaller than in the ATLAS analysis, resulting in larger data and simulation statistical uncertainties. The large difference between the two experiments in the category of jet uncertainties arises because the JES uncertainty in ATLAS is evaluated in different categories mostly using pseudoexperiments, while in CMS the JES uncertainty is introduced as a nuisance parameter in the fit. Further details are discussed in Appendix~\\ref{app:tw}.\nIn the CMS \\mys-channel\\xspace\\ analysis, the uncertainty associated with the size of the simulated samples is evaluated as part of the total statistical uncertainty. More details are discussed in Appendix~\\ref{app:sch}.\n\n\n\\begin{sidewaystable}[!htbp]\n \\centering\n \\caption{Results of the ATLAS and CMS individual $\\ensuremath{\\absvtb^2}\\xspace$ determinations that are the inputs to the overall $\\ensuremath{\\absvtb^2}\\xspace$\n combination together with their experimental uncertainties. The values of $\\ensuremath{\\absvtb^2}\\xspace$ may\n slightly differ from those published for the different analyses since in this paper the theoretical cross-sections used are those marked with $^\\dagger$ in Table~\\ref{tab:xs}. Experimental uncertainties contributing\n less than 1\\% are denoted by $<$0.01. Entries with $-$ mean\n that this uncertainty was not evaluated for this\n analysis. Descriptions of the background categories and of the\n correlations of systematic uncertainties between experiments are presented in\n Appendix~\\ref{app:allUncs}.}\n \\IfPackagesLoaded{adjustbox}{\\begin{adjustbox}{max width=0.80\\textwidth}}{}\n \\begin{tabular}{l|rrrrrrrrr}\n\\hline\n\\hline\n& $t$-channel& $t$-channel& $t$-channel& $t$-channel& $tW$& $tW$& $tW$& $tW$& $s$-channel\\\\\n& ATLAS& CMS& ATLAS& CMS& ATLAS& CMS& ATLAS& CMS& ATLAS\\\\\n& 8 TeV& 8 TeV& 7 TeV& 7 TeV& 8 TeV& 8 TeV& 7 TeV& 7 TeV& 8 TeV\\\\\n\\hline\\hline\n\\T\\B{\\boldmath $\\ensuremath{\\absvtb^2}\\xspace$}&1.06&0.99&1.06&1.05&1.03&1.05&1.07&1.02&0.92\\\\\n\\hline\n\\hline\n\\textbf{Uncertainties:} & & & & & & & & &\\\\\n\\hline\n~~\\textbf{Data statistical}& 0.01& 0.03& 0.03& 0.06& 0.06& 0.09& 0.18& 0.21& 0.15\\\\\n\\hline\n~~\\textbf{Simulation statistical}& 0.01& 0.01& 0.02& 0.02& 0.01& 0.03& 0.02& $-$& 0.11\\\\\n\\hline\n~~\\textbf{Integrated luminosity}& 0.02& 0.03& 0.02& 0.02& 0.05& 0.03& 0.07& 0.04& 0.05\\\\\n\\hline\n~~\\textbf{Theory modelling}& & & & & & & & & \\\\\n~~~~ISR\/FSR, ren.\/fact.\\ scale& 0.04& 0.02& 0.03& 0.04& 0.09& 0.13& 0.05& 0.03& 0.06\\\\\n~~~~NLO match., generator& 0.03& 0.05& 0.02& 0.04& 0.03& $-$& 0.11& $-$& 0.10\\\\\n~~~~Parton shower& 0.02& $-$& $-$& 0.01& 0.02& 0.15& 0.16& 0.10& 0.02\\\\\n~~~~PDF& 0.01& 0.02& 0.03& 0.01& 0.01& 0.02& 0.02& 0.02& 0.03\\\\\n~~~~DS\/DR scheme& $-$& $-$& $-$& $-$& 0.04& 0.02& $-$& 0.06& $-$\\\\\n~~~~Top-quark $p_{\\rm T}$ rew.& $-$& $-$& $-$& $-$& $-$& $<$0.01& $-$& $-$& $-$\\\\\n\\hline\n~~\\textbf{Background normalisation}& & & & & & & & & \\\\\n~~~~Top-quark bkg.& $<$0.01& 0.02& 0.02& 0.01& 0.02& 0.02& 0.06& 0.06& 0.05\\\\\n~~~~Other bkg.\\ from sim.& 0.01& $<$0.01& $<$0.01& 0.03& 0.02& 0.03& 0.09& 0.04& 0.05\\\\\n~~~~Bkg.\\ from data& $<$0.01& 0.02& 0.01& 0.01& $<$0.01& $-$& 0.02& $-$& 0.01\\\\\n\\hline\n~~\\textbf{Jets}& & & & & & & & & \\\\\n~~~~JES common& 0.03& 0.04& 0.08& 0.01& 0.05& 0.04& 0.17& 0.15& 0.05\\\\\n~~~~JES flavour& $<$0.01& $-$& 0.02& $-$& 0.02& $-$& $-$& $-$& 0.01\\\\\n~~~~JetID& $<$0.01& $-$& 0.01& $-$& $<$0.01& $-$& 0.05& $-$& 0.01\\\\\n~~~~JER& $<$0.01& 0.01& 0.02& $<$0.01& 0.07& 0.01& 0.02& 0.04& 0.11\\\\\n\\hline\n~~\\textbf{Detector modelling}& & & & & & & & & \\\\\n~~~~Leptons& 0.02& 0.01& 0.03& 0.04& 0.03& 0.02& 0.07& 0.05& 0.02\\\\\n~~~~HLT (had.\\ part)& $-$& $-$& $-$& 0.02& $-$& $-$& $-$& $-$& $-$\\\\\n~~~~\\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace scale& $<$0.01& $<$0.01& 0.03& $<$0.01& 0.06& $<$0.01& $-$& 0.03& 0.01\\\\\n~~~~\\ensuremath{E_{\\mathrm{T}}^{\\mathrm{miss}}}\\xspace res.& $<$0.01& $-$& $-$& $-$& $<$0.01& $-$& $-$& $-$& 0.01\\\\\n~~~~\\btag& 0.01& 0.02& 0.04& 0.02& 0.01& 0.01& $-$& 0.02& 0.07\\\\\n~~~~Pile-up& $<$0.01& 0.01& $<$0.01& 0.01& 0.03& $<$0.01& 0.11& 0.01& 0.01\\\\\n\\hline\n~~\\textbf{Top-quark mass}& 0.01& $<$0.01& 0.01& $-$& 0.05& 0.05& $-$& $-$& $-$\\\\\n\\hline\n~~\\textbf{Theoretical cross-section}& & & & & & & & & \\\\\n~~~~PDF+$\\alpha_{\\rm s}$& 0.03& 0.03& 0.04& 0.04& 0.06& 0.07& 0.08& 0.07& 0.03\\\\\n~~~~Ren.\/fact.\\ scale& 0.03& 0.03& 0.03& 0.03& 0.03& 0.03& 0.03& 0.03& 0.02\\\\\n~~~~Top-quark mass& 0.01& 0.01& 0.01& 0.01& 0.02& 0.02& 0.02& 0.02& 0.02\\\\\n~~~~$E_{\\text{beam}}$& $<$0.01& $<$0.01& $<$0.01& $<$0.01& $<$0.01& $<$0.01& $<$0.01& $<$0.01& $<$0.01\\\\\n\\hline\n\\hline\n\\textbf{Total systematic uncertainty}&0.09&0.09&0.13&0.10&0.18&0.23&0.34&0.24&0.24\\\\\n\\hline\n\\textbf{Total uncertainty}&0.09&0.10&0.13&0.12&0.19&0.24&0.38&0.32&0.28\\\\\n\\hline\n\\hline\n\\end{tabular}\n \\IfPackagesLoaded{adjustbox}{\\end{adjustbox}}{}\n \\label{tab:uncVtb}\n\\end{sidewaystable}\n\n\\subsection{Results} \n\\label{sec:vtbresult}\nThe combination of $\\ensuremath{\\absvtb^2}\\xspace$ is performed using the inputs from\nall three single-top-quark production modes. Using the same method, the combination of $\\ensuremath{\\absvtb^2}\\xspace$ is also performed separately for each production mode for comparison.\n\nCombining the $\\ensuremath{\\absvtb^2}\\xspace$ values extracted from the \\myt-channel\\xspace\\ and \\ensuremath{tW}\\xspace\\ cross-section measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ from ATLAS and CMS, as well as the ATLAS \\mys-channel\\xspace\\ measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, results in\n\\begin{linenomath}\n \\begin{equation*}\n \\ensuremath{|\\flv\\vtb|^2 = 1.05\\pm0.02\\;(\\text{stat.})\\pm0.06\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.})\\pm0.04\\;(\\text{theo.}) = 1.05\\pm0.08},\n \\label{eq:resultSqd}\n \\end{equation*}\n\\end{linenomath}\nwith a relative uncertainty of $\\ensuremath{7.4\\%}$. The contribution from each experimental uncertainty category to the\ntotal uncertainty in the combined $\\ensuremath{\\absvtb^2}\\xspace$ value is shown in Table~\\ref{tab:vtbRes}. The\ntheory modelling uncertainties in signal and background processes, discussed in\nSection~\\ref{sec:systcat_exp}, dominate the experimental\nuncertainty and the total uncertainty. The theoretical cross-section uncertainty is the\nsecond-largest contribution to the total uncertainty in the combined $\\ensuremath{\\absvtb^2}\\xspace$ value. Changes in the combined $\\ensuremath{\\absvtb^2}\\xspace$ value from using alternative NNLO and NLO+NNLL theoretical predictions for the \\myt-channel\\xspace\\ are less than 1\\%.\n\n\\begin{table}[htbp]\n \\begin{center}\n \\caption{Contributions from each experimental and theoretical\n uncertainty category to the overall $\\ensuremath{\\absvtb^2}\\xspace$ combination. The total uncertainty is\n computed by adding in quadrature all of the individual systematic\n uncertainties (including the integrated luminosity and theoretical\n cross-section) and the statistical uncertainty in data.}\n \\begin{tabular}{l | r | r}\n\\hline\n \\hline\n\\T\\B\n{\\bf Combined} {\\boldmath $\\ensuremath{\\absvtb^2}\\xspace$} & \\multicolumn{2}{c}{1.05} \\\\\n\\hline\n\\hline\n\\multirow{2}{*}{Uncertainty category} & \\multicolumn{2}{c}{Uncertainty} \\\\\\cline{2-3}\n & ~~[\\%] & $\\Delta\\ensuremath{\\absvtb^2}\\xspace$ \\\\\n\\hline\nData statistical & 1.8 & 0.02\\\\\n\\hline\nSimulation statistical & 0.9 & 0.01\\\\\nIntegrated luminosity & 1.3 & 0.01\\\\\nTheory modelling & 4.5 & 0.05\\\\\nBackground normalisation & 1.3 & 0.01\\\\\nJets & 2.6 & 0.03\\\\\nDetector modelling & 1.6 & 0.02\\\\\nTop-quark mass & 0.7 & 0.01\\\\\nTheoretical cross-section & 4.3 & 0.04\\\\\n\\hline\nTotal syst.\\ unc.\\ (excl.\\ lumi.) & 7.1 & 0.07\\\\\nTotal syst.\\ unc.\\ (incl.\\ lumi.) & 7.2 & 0.08\\\\\n\\hline\\hline\n{\\bf Total uncertainty} & 7.4 & 0.08\\\\\n\\hline\n \\hline\n\\end{tabular}\n \\label{tab:vtbRes}\n \\end{center}\n\\end{table}\n\nFigure~\\ref{fig:rhosVtb} illustrates the correlations between the input measurements in the combination. The correlations are all below 0.6. The largest correlations are generally between the measurements in the same experiment at the same centre-of-mass energy, and for those that have large contributions from the same theory modelling components, such as the ATLAS \\mys-channel\\xspace\\ measurement, which has a correlation of over 0.5 with each of the \\ensuremath{tW}\\xspace\\ measurements.\n\nThe BLUE weights for each of the contributing measurements are shown in\nTable~\\ref{tab:wgtVtb}. The \\myt-channel\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ have\nthe largest weight in the combination, followed by the \\myt-channel\\xspace\\ measurements at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}.\nThe $\\ensuremath{tW}\\xspace$ measurements have smaller cross-section uncertainties than the \\mys-channel\\xspace\\ measurements, but, in addition to the correlation between $\\ensuremath{tW}\\xspace$ and \\mys-channel\\xspace\\ measurements, the $\\ensuremath{tW}\\xspace$ measurements are also\nmore correlated with the \\myt-channel\\xspace\\ measurements in each experiment. The negative weights indicate the presence of large correlations between the corresponding measurement and some of the other measurements~\\cite{Valassi:2013bga}.\n\n\\begin{figure}[htbp]\n \\begin{center}\n \\includegraphics[width=0.85\\textwidth]{overallCorrComb}\n \\caption{Correlation matrix of the overall \\ensuremath{\\absvtb^2}\\xspace\\ combination. Each bin\n corresponds to a measurement in a given production mode,\n experiment, and at a given centre-of-mass energy.}\n \\label{fig:rhosVtb}\n \\end{center}\n\\end{figure}\n\n\\begin{table}[htbp]\n \\begin{center}\n \\caption{BLUE weights for the overall $\\ensuremath{\\absvtb^2}\\xspace$ combination.}\n \\begin{tabular}{c|c|c|L{.}}\n\\hline\n\\hline\nProcess & $\\sqrt{s}$ & Experiment & \\multicolumn{1}{c}{BLUE weight} \\\\\n\\hline\n\\multirow{4}{*}{$t$-channel} & \\multirow{2}{*}{8 TeV} & ATLAS&0.56\\\\\n\\cline{3-4}\n& & CMS&0.27\\\\\n\\cline{3-4}\n\\cline{2-3}\n & \\multirow{2}{*}{7 TeV} & ATLAS&0.07\\\\\n\\cline{3-4}\n& & CMS&0.15\\\\\n\\hline\n\\multirow{4}{*}{$tW$} & \\multirow{2}{*}{8 TeV} & ATLAS&0.05\\\\\n\\cline{3-4}\n& & CMS&-0.04\\\\\n\\cline{3-4}\n\\cline{2-3}\n & \\multirow{2}{*}{7 TeV} & ATLAS&-0.02\\\\\n\\cline{3-4}\n& & CMS&0.02\\\\\n\\hline\n$s$-channel & 8 TeV & ATLAS&-0.07\\\\\n\\hline\\hline\n\\end{tabular}\n \\label{tab:wgtVtb}\n \\end{center}\n\\end{table}\n\nThe combined $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ value from the cross-section measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}, including uncertainties in $\\ensuremath{\\sigma_{\\rm theo.}}\\xspace$ for each\nproduction mode, is\n\\begin{linenomath}\n \\begin{equation*}\n \\begin{aligned}\n \\ensuremath{|\\flv\\vtb|=~&~1.02\\pm0.01\\;(\\text{stat.})\\pm0.03\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.})\\pm0.02\\;(\\text{theo.}) \\\\=~&~1.02\\pm0.04\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.}) = 1.02\\pm0.04},\n \\end{aligned}\n \\label{eq:result}\n \\end{equation*}\n\\end{linenomath}\nwith a relative uncertainty of \\ensuremath{3.7\\%}, which improves on the precision of 4.7\\% of the most precise individual $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ extraction, which comes from the ATLAS \\myt-channel\\xspace\\ analysis at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}~\\cite{TOPQ-2015-05}. This is a 30\\% improvement over the Tevatron combination~\\cite{Aaltonen:2015cra}.\n\nThe $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ values are also combined for each production mode, combining across experiments and centre-of-mass energies. For the \\mys-channel\\xspace, the ATLAS and CMS measurements at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}\\ are combined. The results are\n\\begin{linenomath}\n \\begin{equation*}\n \\begin{aligned}\n {t{\\text{-channel}}}: \\ensuremath{|\\flv\\vtb|=~&~1.02\\pm0.01\\;(\\text{stat.})\\pm0.03\\;(\\text{syst.})\\pm0.01\\;(\\text{lumi.})\\pm0.02\\;(\\text{theo.}) \\\\=~&~1.02\\pm0.04\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.}) = 1.02\\pm0.04}, \\\\\n {\\ensuremath{tW}\\xspace}: \\ensuremath{|\\flv\\vtb|=~&~1.02\\pm0.03\\;(\\text{stat.})\\pm0.07\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.})\\pm0.04\\;(\\text{theo.}) \\\\=~&~1.02\\pm0.09\\;(\\text{meas.})\\pm0.04\\;(\\text{theo.}) = 1.02\\pm0.09}, \\\\\n {s{\\text{-channel}}}: \\ensuremath{|\\flv\\vtb|=~&~0.97\\pm0.08\\;(\\text{stat.})\\pm0.12\\;(\\text{syst.})\\pm0.02\\;(\\text{lumi.})\\pm0.02\\;(\\text{theo.}) \\\\=~&~0.97\\pm0.15\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.}) = 0.97\\pm0.15}. \\\\\n \\end{aligned}\n \\end{equation*}\n\\end{linenomath}\nThe relative uncertainties are \\ensuremath{3.9\\%}, \\ensuremath{8.4\\%}\\ and \\ensuremath{15.0\\%}\\ respectively. In all cases, these results\nare more precise than the best individual determinations of $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$, which have uncertainties of\n4.7\\%, 9.9\\% and 20.8\\% for the \\myt-channel\\xspace~\\cite{TOPQ-2015-05},\n\\ensuremath{tW}\\xspace~\\cite{TOPQ-2012-20} and \\mys-channel\\xspace~\\cite{TOPQ-2015-01} analyses respectively.\n\nFigure~\\ref{fig:vtbsum} shows a summary of the $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ combinations. The combination is dominated by the \\myt-channel\\xspace\\ analyses.\n\n\\begin{figure}[!h!tbp]\n \\begin{center}\n \\includegraphics[width=0.82\\textwidth]{figures\/vtb\/vtb_summary.pdf}\n \\caption{The combined $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ value extracted from the \\myt-channel\\xspace\\ and \\ensuremath{tW}\\xspace\\ cross-section measurements at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ from ATLAS and CMS, as well as the ATLAS \\mys-channel\\xspace measurement at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}, is shown together with the combined $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ values for each production mode.\n The theoretical predictions for \\myt-channel\\xspace\\ and \\mys-channel\\xspace\\ production are computed at NLO accuracy, while the theoretical predictions for $\\ensuremath{tW}\\xspace$ are calculated at NLO+NNLL accuracy. The $\\ensuremath{\\sigma_{\\rm theo.}}\\xspace$ uncertainties used to compute $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ include scale, PDF+$\\alpha_{\\rm s}$, $m_{\\ensuremath{t}\\xspace}$, and $E_{\\rm beam}$ variations.\n }\n \\label{fig:vtbsum}\n\\end{center}\n\\end{figure}\n\n\\subsection{Stability tests}\n\\label{sec:tests}\n\nThe stability of the combination of the $\\ensuremath{\\absvtb^2}\\xspace$ values to variations in the correlation assumptions, discussed in Section~\\ref{sec:systcat}, is checked for the dominant uncertainty contributions. The correlation values are varied for the theory modelling, JES, and the most important contributions to the theoretical cross-section predictions (i.e.\\ PDF+$\\alpha_{\\rm s}$ and scale). Because of the scheme that is used for the correlations, stability tests are also performed for the uncertainties associated with the integrated luminosity. Figure~\\ref{fig:vtbstabilities} summarises the results of these stability tests, where the correlations between ATLAS and CMS (and also between centre-of-mass energies for the integrated luminosity) are varied.\n\n\\begin{figure}[!h!tbp]\n \\begin{center}\n \\includegraphics[width=1.0\\textwidth]{stability\/stability_style2.pdf}\n \\caption{Results of the stability tests performed by varying\n of the correlation assumptions in different uncertainty\n categories: theory modelling (scales and radiation modelling,\n NLO matching, and PS and hadronisation), JES, dominant\n theoretical cross-section predictions (i.e.\\ PDF+$\\alpha_{\\rm s}$\n and scale) and integrated luminosity. Two or three variations are considered depending on the\n uncertainty category. The corresponding relative shifts\n (with shift = varied $-$ nominal) in the central value, $\\Delta\\ensuremath{\\absvtb^2}\\xspace\/\\ensuremath{\\absvtb^2}\\xspace$, and in its uncertainty,\n $\\Delta(\\delta\\ensuremath{\\absvtb^2}\\xspace)\/(\\delta\\ensuremath{\\absvtb^2}\\xspace)$, are shown.}\n \\label{fig:vtbstabilities}\n \\end{center}\n\\end{figure}\n\nThe uncertainties in the theory modelling category (i.e.\\ scales and radiation modelling, NLO matching, and PS and\nhadronisation) are varied from their default value of fully\ncorrelated to half correlated and to the more extreme tests of uncorrelated and half anti-correlated. The JES\ncategory is varied from its default value of uncorrelated to half correlated and half anti-correlated and the more extreme variation of fully correlated. The theoretical cross-section uncertainties, PDF+$\\alpha_{\\rm s}$ and\nscale, are varied from their default values of fully correlated to\nhalf correlated, uncorrelated and half anti-correlated. For the integrated luminosity, the correlation between ATLAS and CMS is varied from its default value of 30\\% correlated to half correlated and uncorrelated. The correlation between different centre-of-mass energies for each experiment is varied from the default of uncorrelated to half and fully\ncorrelated. The correlation of the theoretical scale uncertainty between different processes is also tested. For all variations, the relative changes in the central value of the combined $\\ensuremath{|f_{\\rm LV}\\vtb|}\\xspace$ are significantly smaller ($<$$0.5\\%$) than the relative total uncertainty of \\ensuremath{3.7\\%}. Additionally, the relative changes in the total uncertainty are below 0.004, i.e., less than $10\\%$ of the total uncertainty of \\ensuremath{0.04}. \nThese tests show that the result of the combination is robust and does not critically depend on any of the correlation assumptions. The cross-section combinations similarly do not depend significantly on any of the correlation assumptions.\n \n\\IfPackagesLoaded{adjustbox}{\\FloatBarrier}{}\n\\section{Summary}\n\\label{sec:sum}\nThe combinations of single-top-quark production cross-section measurements in the \\myt-channel\\xspace, \\ensuremath{tW}\\xspace, and \\mys-channel\\xspace\\ production modes are presented, using data from LHC $pp$ collisions collected by the ATLAS and CMS Collaborations. The combinations for each production mode are performed at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}, using data corresponding to integrated luminosities of 1.17 to 5.1~\\ensuremath{\\rm fb^{-1}}\\xspace at \\ensuremath{\\sqrt{s}=\\lhcenergySeven}, and of 12.2 to 20.3~\\ensuremath{\\rm fb^{-1}}\\xspace at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}. \nThe combined \\myt-channel\\xspace\\ cross-sections are found to be $\\ensuremath{67.5\\pm5.7\\,\\,\\text{pb}}$ and $\\ensuremath{87.7\\pm5.8\\,\\,\\text{pb}}$ at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ respectively. The values of the combined \\ensuremath{tW}\\xspace\\ cross-sections at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ are $\\ensuremath{16.3\\pm4.1\\,\\,\\text{pb}}$ and $\\ensuremath{23.1\\pm3.6\\,\\,\\text{pb}}$ respectively. For the \\mys-channel\\xspace\\ cross-section, the combination yields $\\ensuremath{4.9\\pm1.4\\,\\,\\text{pb}}$ at \\ensuremath{\\sqrt{s}=\\lhcenergyEight}. \nThe square of the magnitude of the CKM matrix element $\\ensuremath{V_{\\myt\\myb}}\\xspace$ multiplied by a form factor accounting for possible contributions from physics beyond the SM, \\ensuremath{f_{\\rm LV}}\\xspace, is determined from each production mode at each centre-of-mass energy, using the ratio of the measured cross-section to its theoretical prediction, and assuming that the top-quark-related CKM matrix elements obey the relation $|\\ensuremath{V_{td}}\\xspace|,|\\ensuremath{V_{ts}}\\xspace| \\ll |\\ensuremath{V_{\\myt\\myb}}\\xspace|$. The values of $\\ensuremath{\\absvtb^2}\\xspace$ extracted from individual ratios at \\ensuremath{\\sqrt{s}=7 \\text{ and } \\lhcenergyEight}\\ yield a combined value of \\ensuremath{|\\flv\\vtb| = 1.02\\pm0.04\\;(\\text{meas.})\\pm0.02\\;(\\text{theo.})}. All combined measurements are consistent with their corresponding SM predictions.\n \n\n\n\\clearpage\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nElectronic transport in mesoscopic systems is usually affected by the\nCoulomb interactions between the electrons. The interactions manifest\nthemselves most dramatically in small tunnel junctions. If a small\nmetallic grain is connected to a lead by a tunnel junction, an electron\ntunneling into such a grain charges it by the elementary charge $e$ and\nincreases the energy of the system by $e^2\/2C$. For a small grain, the\ntypical capacitance of the system can be very small, so that the\ncharging energy $e^2\/2C$ can significantly exceed the temperature of the\nsystem $T$. In this regime one observes the Coulomb blockade: strong\nsuppression of the tunneling into the grain at low temperatures. Various\naspects of the Coulomb blockade have been intensively studied in recent\nyears, both experimentally and theoretically.\\cite{Averin,Grabert}\n\nRecently it has become possible\\cite{Kastner} to observe the Coulomb\nblockade in semiconductor heterostructures. Instead of a metallic grain,\nin such experiments one creates a small region of the two-dimensional\nelectron gas (2DEG)---a quantum dot---separated from the rest of the\n2DEG by potential barriers created by applying negative voltage to the\ngates, Fig.~\\ref{fig:1}. Experimentally, one measures the linear\nconductance through the quantum dot as a function of the voltage $V_g$\napplied to the central gates. The role of the gate voltage is to change\nthe electrostatic energy of the system,\n\\begin{equation}\nE = \\frac{(Q-eN)^2}{2C},\n \\label{electrostatic_energy}\n\\end{equation}\nwhere $Q$ is the charge of the quantum dot, and $N$ is a dimensionless\nparameter proportional to the gate voltage and the capacitance $C_g$\nbetween the dot and the gate, $N=C_g V_g\/e$. If the potential barriers\nseparating the quantum dot from the leads are high, the charge $Q$ is\ninteger in units of $e$. Thus, the ground state of the system corresponds\nto some integer value of charge, and the states with different values of\n$Q$ are separated from it by the electrostatic gap $\\Delta\\sim e^2\/2C$. It\nis important to note that at the values of the gate voltage corresponding\nto half-integer values of $N$ the gap vanishes, i.e., the ground state is\ndegenerate in charge $Q$. For example, at $N=\\frac{1}{2}$ the ground state\ncan have either charge $Q=0$ or $e$. Therefore at half-integer $N$ the\ntunneling of an electron into or out of the quantum dot does not lead to\nthe increase of the electrostatic energy of the system, and the Coulomb\nblockade is lifted. As a result the conductance of the structure in\nFig.~\\ref{fig:1} shows periodic peaks as a function of $N$. The Coulomb\nblockade peaks in linear conductance as a function of the gate voltage are\nreadily observed in the experiments (see, e.g.,\nRef.~\\onlinecite{Kastner}).\n\n\n\\vfill\n\\begin{figure}\n\\narrowtext\n\\epsfxsize=3.1in\n\\hspace{.2em}\n\\epsffile{4chan1.ps}\n\\vspace{1.5\\baselineskip}\n\\caption{Schematic view of a quantum dot connected to two bulk 2D\n electrodes. The dot is formed by applying negative voltage to the gates\n (shaded). Solid line shows the boundary of the 2D electron gas (2DEG).\n Electrostatic conditions in the dot are controlled by the voltage\n applied to the central gates. Voltage $V_{L,R}$ applied to the auxiliary\n gates controls the transmission probability through the left and right\n constrictions.}\n\\label{fig:1}\n\\end{figure}\n\\vspace{\\baselineskip}\n\n\n\n\nThe shapes and heights of the peaks in conductance are usually described\nwithin the framework of the orthodox theory of the Coulomb\nblockade.\\cite{Averin} In this approach the conductance is found from the\nbalance of the rates of tunneling of electrons between the leads and the\ndot, with the rates being calculated in the lowest order perturbation\ntheory in the tunneling matrix elements. The resulting linear conductance\nhas the form:\\cite{Glazman}\n\\begin{equation}\n G = \\frac{1}{2}\\frac{G_LG_R}{G_L + G_R}\n \\frac{\\Delta\/T}{\\sinh(\\Delta\/T)},\n \\quad \\Delta=\\frac{e^2}{C}\\delta N.\n \\label{orthodox}\n\\end{equation}\nHere $G_L$ and $G_R$ are the conductances of the left and right barriers,\n$\\delta N$ is the deviation of $N$ from the nearest half-integer value.\nAway from the center of the peak, the tunneling of an electron into the\nquantum dot is suppressed, because only a small fraction of electrons has\nthe energy sufficient to overcome the electrostatic gap $\\Delta$. Thus the\nconductance (\\ref{orthodox}) decays exponentially at $T\\to0$. At very low\ntemperatures, however, another mechanism---usually referred to as {\\it\n inelastic co-tunneling}---dominates the conductance. Instead of a\nreal process of tunneling from a lead to the dot, electron tunnels from\nthe left lead to a virtual state in the dot, and then another electron\ntunnels from the dot to the right lead. As a result of such a two-step\nprocess, one electron is transferred from the left lead to the right one,\nand the charge of the quantum dot is unchanged. The corresponding\nconductance has only a power-law dependence on temperature,\\cite{AN}\n\\begin{equation}\nG=\\frac{\\pi\\hbar G_LG_R}{3e^2}\\left(\\frac{T}{\\Delta}\\right)^2,\n\\quad\nT\\ll\\Delta.\n\\label{cotunneling}\n\\end{equation}\nIn the case of weak tunneling, $G_L, G_R\\ll e^2\/\\hbar$, the results\n(\\ref{orthodox}) and (\\ref{cotunneling}) give the full description of the\nCoulomb blockade peaks in linear conductance. Near the center of the peak\nthe conductance is given by Eq.~(\\ref{orthodox}), whereas the tails of the\npeaks are dominated by the co-tunneling contribution (\\ref{cotunneling}).\n\nIn a number of recent experiments\\cite{van,Pasquier,Waugh,Molenkamp} the\ntunneling through quantum dots was studied in the regime when one or more\nquantum point contacts were tuned to the regime of strong tunneling,\ncorresponding to conductance $G\\to e^2\/\\pi\\hbar$. In particular, it was\ndemonstrated by van der Vaart {\\it et al.\\\/}\\cite{van} that the Coulomb\nblockade oscillations of linear conductance persist as long as the\nconductances of both barriers are below the conductance quantum\n$e^2\/\\pi\\hbar$. Most of the existing theoretical works on the\nstrong-tunneling regime of the Coulomb\nblockade\\cite{GS,PZ,Falci,Schoeller} discuss the case of metallic\nsystems, where one can achieve large conductance of a tunnel junction\nbetween the grain and a lead by increasing the number of transverse modes\nwhile keeping the transmission coefficient ${\\cal T}$ small. Such theories\ncannot be applied to the experiments in semiconductor heterostructures,\nwhere the point contacts typically allow only a single transverse mode,\nbut the transmission coefficient can be tuned to ${\\cal T} \\to 1$. The\ntheoretical works for this case either do not concentrate on the Coulomb\nblockade oscillations\\cite{Flensberg} or study the oscillations of\nequilibrium properties of the system,\\cite{Matveev1} rather than its\nconductance.\n\nIn this paper we develop a theory of the Coulomb blockade oscillations of\nconductance in the regime of strong tunneling. It was shown in\nRef.~\\onlinecite{Matveev1} that the average charge of a quantum dot\nconnected to a lead by a single-mode quantum point contact shows small\nperiodic oscillations as a function of the gate voltage as long as the\ntransmission coefficient ${\\cal T}$ is below unity. Similarly, we will\nshow that the conductance through a quantum dot shows small periodic\noscillations at ${\\cal T} < 1$. One should note, however, that unlike the\naverage charge,\\cite{Matveev1} the amplitude of the oscillations of conductance\ndepends\non temperature and grows as $T$ is lowered. In the limit $T\\to 0$, small\noscillations develop into sharp periodic peaks, and the off-peak\nconductance vanishes as $T^2$, thus reproducing the characteristic\nquadratic temperature dependence for inelastic co-tunneling\n(\\ref{cotunneling}).\n\nIn section \\ref{weak} we use the analogy\\cite{GM,Matveev2} between the\nCoulomb blockade and multichannel Kondo problem to discuss the peak value\nof conductance at $T\\to0$. This discussion is done for the weak-tunneling\ncase, but the conclusion that the peak conductance is of the order of\n$e^2\/\\hbar$ is valid for the strong tunneling as well. The theory of\nconductance oscillations in the strong-tunneling regime is developed in\nsections \\ref{spinless} and \\ref{withspin}. In section \\ref{scaling} we\ndiscuss the scaling properties of our problem. In particular, we show that\nthe quadratic temperature dependence of the off-peak conductance at\n$T\\to0$ is universal and holds for arbitrary transmission coefficients.\n\n\\section{Conductance peaks in the weak-tunneling case}\n\\label{weak}\n\nLet us first consider the weak-tunneling case. The linear conductance is\ngiven by Eqs.~(\\ref{orthodox}) and (\\ref{cotunneling}). One should note\nthat in the derivation of Eqs.~(\\ref{orthodox}) and (\\ref{cotunneling})\nonly the lowest-order terms of the perturbation theory in tunneling\namplitudes were taken into account. This is usually justified at $G_L, G_R\n\\ll e^2\/\\hbar$. However, it was shown in Ref.~\\onlinecite{GM} that near a\nhalf-integer $N$ and at low temperatures the higher-order corrections lead\nto large logarithmic renormalizations of the tunneling amplitudes. The\nnature of these renormalizations can be understood in terms of the\nanalogy\\cite{Matveev2} between the Coulomb blockade and the Kondo problem.\nWe will now discuss this analogy and correct the results (\\ref{orthodox})\nand (\\ref{cotunneling}) to take into account the renormalizations of the\ntunneling amplitudes.\n\nFor simplicity we start with the quantum dot weakly coupled to a {\\it\n single\\\/} lead. Such a system can be described by the following\nHamiltonian:\n\\begin{eqnarray}\nH &=& \\sum_{k} \\epsilon_k a_k^\\dagger a_k^{}\n + \\sum_{p} \\epsilon_p a_p^\\dagger a_p^{}\n + E_C (\\hat n - N)^2\n\\nonumber\\\\\n & & +\\sum_{kp} (v_t^{} a_k^\\dagger a_p^{} + v_t^* a_p^\\dagger a_k^{}).\n \\label{tunnel_hamiltonian}\n\\end{eqnarray}\nHere $a_k$ and $a_p$ are the annihilation operators for the electrons in\nthe lead and in the dot respectively, and $E_C$ is the characteristic\ncharging energy, $E_C = e^2\/2C$; coupling between the dot and the lead is\ndescribed by tunneling Hamiltonian with the matrix element $v_t^{}$.\nFinally, $\\hat n$ is the operator of the number of electrons in the dot,\ncounted from that at the ground state of the system without tunneling,\n$\\hat n = \\sum (a_p^\\dagger a_p^{} - \\langle a_p^\\dagger a_p^{}\n\\rangle_0^{})$.\n\nDue to the last term in Eq.~(\\ref{tunnel_hamiltonian}) electrons can\ntunnel from the lead to a virtual state in the dot. The energy of such a\nvirtual state is of the order of $E_C$, and to complete this scattering\nprocess an electron has to tunnel from the dot to the lead. It is\ninstructive to calculate the second-order amplitude of scattering between\ntwo states in the lead, $k$ and $k'$, near the Fermi level:\n\\begin{eqnarray}\n T_{k\\to k'} = \\sum_p && \\left[\n \\frac{|v_t^{}|^2 \\theta(\\epsilon_p)}\n {E_C N^2 - E_C(1-N)^2 - \\epsilon_p}\n\\right.\\nonumber\\\\\n && \\left.\n - \\frac{|v_t^{}|^2 \\theta(-\\epsilon_p)}\n {E_C N^2 - E_C(1+N)^2 + \\epsilon_p}\n \\right].\n \\label{2_order_sum}\n\\end{eqnarray}\nHere the second term corresponds to the process in which first an electron\nfrom the filled state $p$ in the dot tunnels to the state $k'$ in the\nlead, and then the electron $k$ fills the hole in the dot. Clearly, the\ntotal amplitude diverges logarithmically at $N\\to\\pm\\frac12$,\n\\begin{equation}\n T_{k\\to k'} = -\\nu_d |v_t^{}|^2 \\ln\\frac{1+2N}{1-2N},\n \\label{2_order_log}\n\\end{equation}\nwhere $\\nu_d$ is the density of states in the dot. The logarithmic\nsingularity at $N\\to\\frac12$ originates from the degeneracy of the\nelectrostatic energy for states with $n=0$ and $n=1$ extra electrons in\nthe dot.\n\nLet us now concentrate on the case of the dimensionless gate voltage $N$\nvery close to $\\frac{1}{2}$. In this case the logarithmic renormalizations\nare large, and to find the scaling properties of the scattering\namplitudes, we can restrict our consideration to the energy scales smaller\nthan $E_C$. Thus when constructing the perturbation series for a tunneling\namplitude we will neglect all the processes in which the system visits\nvirtual states with the charge $n$ different from 0 or 1. Another\nimportant property of the Hamiltonian (\\ref{tunnel_hamiltonian}) is that\nthe perturbation---the tunneling term---always changes the charge of the\ndot by 1. Suppose that at the first step of a high-order tunneling process\nan electron tunnels from the lead to the dot thus changing $n$ from 0 to\n1. Then at the next step an electron will have to tunnel from the dot to\nthe lead (changing $n=1\\to0$), because otherwise a virtual state with\n$n=2$ and correspondingly large energy $\\sim E_C$ would have been created.\nAt the third step an electron tunnels from the lead to the dot, changing\n$n=0\\to1$, and so on.\n\nOne can easily see that exactly the same perturbation series for\nscattering amplitudes is obtained from the Hamiltonian of the anisotropic\nKondo model,\n\\begin{eqnarray}\nH &=& \\sum_{k\\alpha} \\epsilon_k a_{k\\alpha}^\\dagger a_{k\\alpha}^{}\n + h S^z\n\\nonumber\\\\\n & & +J_\\perp\\! \\sum_{kk'\\alpha\\alpha'}\\!\n (\\sigma^{+}_{\\alpha\\alpha'} S^{-} + \\sigma^{-}_{\\alpha\\alpha'} S^{+})\n \\,a_{k\\alpha}^\\dagger a_{k'\\alpha'}^{},\n \\label{Kondo}\n\\end{eqnarray}\nwhere $\\sigma^{\\pm}$ are the linear combinations of Pauli matrices,\n$\\sigma^{\\pm}= \\sigma^x \\pm \\sigma^y$, and {\\boldmath $S$} is the impurity\nspin, $S^{\\pm}=S^x\\pm S^y$. A formal discussion of this mapping can be\nfound in Ref.~\\onlinecite{Matveev2}; here we only mention the relations\nbetween the parameters of the Hamiltonians (\\ref{tunnel_hamiltonian}) and\n(\\ref{Kondo}). The values of the electron spin $\\alpha=\\ \\uparrow$ and\n$\\downarrow$ correspond to the electron being in the lead and the dot\nrespectively. Similarly, the values of the impurity spin $S^z=\\\n\\downarrow$ and $\\uparrow$ correspond to the states with the number of\nparticles in the dot $n=0$ and 1. Magnetic field $h$ is identified with\nthe energy splitting of the $n=0$ and $n=1$ states, $h=2E_C(\\frac12-N)$.\nFinally, the tunneling matrix element $v_t^{}$ is mapped to the exchange\nconstant $J_\\perp$.\n\nIt is well known that at low temperatures and magnetic fields the exchange\nconstant $J$ of the Kondo model renormalizes logarithmically. As a result\n$J$ grows, and the system approaches the strong-coupling fixed point. In\nthe context of the Coulomb blockade problem, this means that at\n$N\\to\\frac12$ and $T\\to0$ the tunneling amplitudes grow until the\nconductance of the barrier becomes of order $e^2\/\\hbar$.\n\nUsing the analogy with the Kondo problem, one can actually find the\nrenormalizations of the tunneling amplitudes. In the leading logarithm\napproximation one can use the poor man's scaling technique.\\cite{Anderson}\nThe result\\cite{GM} can be written down in terms of the renormalization of\nthe transmission coefficient ${\\cal T}$ of the barrier between the lead\nand the dot,\n\\begin{equation}\n{\\cal T} = \\frac{{\\cal T}_0}\n {\\cos^2\\bigl(\\frac{1}{\\pi}\\sqrt{{\\cal T}_0}\\xi\\bigr)},\n\\quad\n\\xi = \\ln\\frac{E_C}{\\max\\{T,\\Delta\\}}.\n \\label{poor_man}\n\\end{equation}\nHere ${\\cal T}_0$ is the bare value of the transmission coefficient, and\n$\\Delta = E_C|1-2N|$. When ${\\cal T}$ becomes of the order of unity, the\nleading logarithm approximation fails. Thus the result (\\ref{poor_man}) is\nvalid as long as the renormalized transmission coefficient is small,\n${\\cal T}\\ll 1$.\n\nIn this regime one can still use the weak-tunneling results\n(\\ref{orthodox}) and (\\ref{cotunneling}) for the conductance through a dot\nconnected to two leads, taking into account the renormalizations of the\nconductances of the barriers. This can be done by noting the relation\nbetween the conductance of a barrier $G_b$ and its transmission\ncoefficient,\n\\begin{equation}\nG_b=\\frac{e^2}{2\\pi\\hbar}\\times\n \\begin{cases}{\n {\\cal T}, & without spins,\\cr\n 2{\\cal T}, & with spins.\n }\\end{cases}\n \\label{barrier_conductance}\n\\end{equation}\nIt is worth mentioning that within the leading logarithm approximation the\ntransmission coefficients of two barriers renormalize independently, so\none can use Eq.~(\\ref{poor_man}) for both left and right barriers.\n\nLet us find the temperature dependence of the peak value of the\nconductance in the case of symmetric barriers, i.e., assume $N=\\frac12$\nand $G_L = G_R = G_b = (e^2\/\\pi\\hbar) {\\cal T}$. According to\nEq.~(\\ref{orthodox}), the peak conductance is $\\frac14 G_b$. Taking into\naccount the renormalization (\\ref{poor_man}) of the transmission\ncoefficient, we find\n\\begin{equation}\nG = \\frac{e^2}{4\\pi\\hbar}\n \\frac{{\\cal T}_0}\n {\\cos^2\\left[\\frac{1}{\\pi}\\sqrt{{\\cal T}_0}\\ln(E_C\/T)\\right]}.\n \\label{peak_conductance}\n\\end{equation}\nAs expected, when the temperature is lowered, the conductance grows. This\nresult is valid only at temperatures exceeding the characteristic Kondo\ntemperature of this problem, $T_K\\simeq E_C\\exp(-\\pi^2\/2\\sqrt{{\\cal\n T}_0})$. At too low temperatures, $T\\lesssim T_K$, the conductance is\nof the order of $e^2\/\\hbar$, transmission coefficient ${\\cal T} \\sim 1$,\nand the result of the leading logarithm approximation\n(\\ref{peak_conductance}) is no longer applicable.\n\nIt is also interesting to study the tails of the conductance peaks. If the\ngate voltage $N$ is not precisely half-integer, then at low temperatures,\n$T\\ll\\Delta=2E_C|N-\\frac12|$, one can still use the inelastic co-tunneling\nresult (\\ref{cotunneling}) with the values of conductances of the barriers\nrenormalized according to Eqs.~(\\ref{poor_man}) and\n(\\ref{barrier_conductance}). The resulting conductance has the form:\n\\begin{equation}\nG=\\frac{e^2{\\cal T}_L{\\cal T}_R}{3\\pi\\hbar}\n \\frac{T^2}\n {\\Delta^2\n \\cos^2\\left[\\frac{1}{\\pi}\\sqrt{{\\cal T}_L}\\ln\\frac{E_C}{\\Delta}\\right]\n \\cos^2\\left[\\frac{1}{\\pi}\\sqrt{{\\cal T}_R}\\ln\\frac{E_C}{\\Delta}\\right]}.\n \\label{tails}\n\\end{equation}\nSimilarly to Eq.~(\\ref{peak_conductance}), the result (\\ref{tails}) is\nvalid only as long as ${\\cal T}\\ll 1$. In this case the condition allows\nonly the values of the gate voltage not too close to the resonance,\n$\\Delta\\gg T_K$.\n\nAs we have seen, if one starts with the weak-tunneling case, in the\nvicinity of the resonances in $G(V_g)$ the tunneling amplitude grows, and\nthe system approaches the strong-coupling regime. Therefore one can use\nthe results of the weak-tunneling approach to find the universal features\nof the peaks in conductance which should also persist in the\nstrong-tunneling case. Such features are: (i) the $T\\to0$ value of the\npeak conductance (\\ref{peak_conductance}) for symmetric barriers is of the\norder of $e^2\/\\hbar$, and (ii) away from the centers of peaks the\ntemperature dependence of the conductance at $T\\to0$ is $G\\propto T^2$. It\nis worth noting that unlike the temperature dependence of the co-tunneling\nconductance (\\ref{cotunneling}), its $1\/\\Delta^2$ dependence of the gate\nvoltage is not universal, see Eq.~(\\ref{tails}).\n\n\n\\section{Strong co-tunneling of spinless electrons}\n\\label{spinless}\n\nWe will start our treatment of the strong-tunneling case with a simplified\nmodel of spinless electrons. As we will see, in this case a complete\nsolution of the problem is possible. Experimentally, the spinless case can\nbe achieved by applying an in-plane magnetic field, sufficiently strong to\npolarize the electron gas in the constrictions.\n\n\\subsection{Theoretical model}\n\nWe assume that each of the quantum point contacts in Fig.~\\ref{fig:1}\nallows the transmission of electrons in only one mode. The transmission\ncoefficients of the two contacts ${\\cal T}_L$ and ${\\cal T}_R$ are assumed\nto be close to unity. This enables us to start with the Hamiltonian\ncorresponding to perfect transmission, ${\\cal T}_L = {\\cal T}_R = 1$, and\nthen construct the perturbation theory in small reflection amplitudes,\n$r_L$ and $r_R$. Following Ref.~\\onlinecite{Flensberg}, we will describe\neach of the quantum point contacts by a system of one-dimensional (1D)\nelectrons. The formal derivation of this model\\cite{Matveev1} is based on\nthe fact that if the boundaries of the channel are very smooth, the\ntransport though it can be described in terms of adiabatic\nwavefunctions\\cite{Lesovik} having essentially a 1D form.\n\nAnother simplification of the problem is due to the fact that the typical\nenergies involved in the Coulomb blockade phenomenon are of the order of\nthe charging energy $E_C$, which is much smaller than the Fermi energy.\nThis allows us to linearize the spectra of the 1D electrons near the\nFermi points and to formulate the model in terms of the left- and\nright-moving fermions:\n\\begin{eqnarray}\nH_0\\! =\\! i\\hbar v_F\\!\\int\\!\\! &&\n \\left[\n \\psi^\\dagger_{L1}(x)\\nabla\\psi_{L1}(x)\n \\!-\\!\\psi^\\dagger_{L2}(x)\\nabla\\psi_{L2}(x)\n \\right.\n\\nonumber\\\\\n &&\\left.\\!\\!\\!\n +\\psi^\\dagger_{R1}(x)\\nabla\\psi_{R1}(x)\n \\!-\\!\\psi^\\dagger_{R2}(x)\\nabla\\psi_{R2}(x)\n \\right] \\!dx.\n \\label{H_0_fermionic}\n\\end{eqnarray}\nHere $\\psi_{L1}$ and $\\psi_{L2}$ are the annihilation operators for the\nelectrons of the left contact moving to the left and right respectively;\n$\\psi_{R1}$ and $\\psi_{R2}$ describe the left- and right-moving electrons\nof the right contact; $\\nabla\\equiv\\frac{\\partial}{\\partial x}$. We\nassociate the centers of the constrictions with points $x=0$, meaning that\nelectrons in the dot are described by $\\psi_{L1(2)}$ at $x>0$ and by\n$\\psi_{R1(2)}$ at $x<0$. Thus the charging energy can be presented as\n\\begin{equation}\n H_C = E_C(\\hat n - N)^2,\n\\label{H_C}\n\\end{equation}\nwhere\n\\[\n\\hat n = \\int_{0}^{\\infty}\\!\\sum_{j=1,2}\n \\!\\left[:\\!\\psi^\\dagger_{Lj}(x)\\psi_{Lj}^{}(x)\\!\\!:\n +:\\!\\psi^\\dagger_{Rj}(-x)\\psi_{Rj}^{}(-x)\\!\\!:\\right]\\!dx.\n\\]\nHere $:\\!\\hat X\\!\\!:$ denotes the normal-ordered operator $\\hat X$.\nFinally, the presence of weak backscattering in the contacts, described by\nsmall reflection amplitudes $|r_L|, |r_R| \\ll 1$, gives rise to the term\nmixing the left- and right-moving particles,\n\\begin{equation}\nH' = \\hbar v_F\\!\\left[\n |r_L|\\psi^\\dagger_{L1}(0)\\psi_{L2}^{}(0)\n +|r_R|\\psi^\\dagger_{R1}(0)\\psi_{R2}^{}(0)\n +\\mbox{h.c.}\\right]\\!.\n \\label{backscattering_fermionic}\n\\end{equation}\n\nThe Hamiltonian $H=H_0+H_C+H'$ gives the complete description of the\ntransport through a quantum dot connected to two leads. In our model the\ntwo contacts are described by two independent 1D systems. Thus we neglect\nthe possibility of coherent transport of electrons from one quantum point\ncontact to the other. Such processes lead to the additional mechanism of\ntransport through the dot called {\\it elastic co-tunneling.\\\/} It is\nknown,\\cite{AN} however, that at temperatures exceeding the level spacing\n$\\varepsilon$ in the dot the elastic co-tunneling contribution to the\ntotal conductance is much smaller than that of inelastic co-tunneling.\nThus we will concentrate on the case $T\\gg\\varepsilon$ which is easy to\nsatisfy in the experiments, and in the following calculations we will\nassume $\\varepsilon=0$. At $T\\lesssim\\varepsilon$ the conductance becomes\ntemperature-independent and can be estimated by replacing $T$ by\n$\\varepsilon$ in our results.\n\nThe solution of the problem with the Hamiltonian given by\nEqs.~(\\ref{H_0_fermionic})--(\\ref{backscattering_fermionic}) is not\ntrivial because the interaction term $H_C$ is not quadratic in fermion\noperators. To overcome this difficulty, we will use the bosonization\ntechnique\\cite{Sol} which enables one to treat the interaction in 1D\nsystems exactly. Following Ref.~\\onlinecite{Sol}, we present the fermion\noperators as\n\\begin{equation}\n \\begin{array}{rcl}\n \\psi_{L1}(x) &=& \\left(\\frac{D}{2\\pi\\hbar v_F}\\right)^{1\/2}\n \\eta_{L1}^{} e^{i\\varphi_{L1}^{}(x)},\\\\\n \\psi_{L2}(x) &=& \\left(\\frac{D}{2\\pi\\hbar v_F}\\right)^{1\/2}\n \\eta_{L2}^{} e^{i\\varphi_{L2}^{}(x)}.\n \\end{array}\n \\label{bosonization}\n\\end{equation}\nHere $D$ is the bandwidth, and the two independent bosonic fields\n$\\varphi_{L1}^{}$ and $\\varphi_{L2}^{}$ satisfy the following commutation\nrelations:\n\\begin{equation}\n[\\varphi_{L1}^{}(x), \\varphi_{L1}^{}(y)]\n = -[\\varphi_{L2}^{}(x), \\varphi_{L2}^{}(y)]\n = -i\\pi{\\rm sgn}(x-y).\n \\label{commutators}\n\\end{equation}\nGiven (\\ref{commutators}), the operators (\\ref{bosonization}) satisfy\nusual anticommutation relations. To ensure proper anticommutation\nrelations for fermions at different branches, we have added\nthe local Majorana fermions $\\eta_{L1}^{}$ and $\\eta_{L2}^{}$.\n\n\nThe electrons of the right contact are\nbosonized in a similar way. One can then rewrite the Hamiltonian\n(\\ref{H_0_fermionic})--(\\ref{backscattering_fermionic}) in terms of\n$\\varphi_L$ and $\\varphi_R$. It will be more convenient, however, to\npresent the Hamiltonian in terms of slightly different bosonic fields:\n\\begin{equation}\n\\phi_L = \\frac{1}{2}(\\varphi_{L2} - \\varphi_{L1}),\n\\quad\n\\Pi_L = -\\frac{1}{2\\pi} \\nabla (\\varphi_{L1} + \\varphi_{L2}).\n \\label{good_fields}\n\\end{equation}\nThe new variables satisfy the commutation relations typical for the\ndisplacement of a 1D elastic medium and its momentum density,\n\\begin{equation}\n[\\phi_L(x), \\Pi_L(y)] = i \\delta(x-y).\n \\label{good_commutator}\n\\end{equation}\nThe form of the free-electron Hamiltonian (\\ref{H_0_fermionic}) supports\nthis observation,\n\\begin{eqnarray}\nH_0 = \\frac{\\hbar v_F}{2}\\int\n &&\\left\\{\\frac{1}{\\pi}[\\nabla \\phi_L(x)]^2 + \\pi \\Pi_L^2(x)\\right.\n\\nonumber\\\\\n &&\\hspace{.7em}\n \\left. +\\frac{1}{\\pi}[\\nabla \\phi_R(x)]^2 + \\pi \\Pi_R^2(x)\n \\right\\}dx.\n \\label{H_0_bosonized}\n\\end{eqnarray}\n\nThis qualitative interpretation allows one to predict the correct\nbosonized form of the charging energy (\\ref{H_C}). One should expect that\nthe charge of the dot is proportional to the difference of the charge\n$Q_L\\sim e\\phi_L(0)$ brought into the dot through the left contact and the\ncharge $Q_R\\sim e\\phi_R(0)$ carried away through the right contact.\nExplicit calculation gives\n\\begin{equation}\nH_C = \\frac{E_C}{\\pi^2}\n \\left[\\phi_R(0) - \\phi_L(0) -\\pi N\\right]^2.\n \\label{H_C_bosonized}\n\\end{equation}\nIt is important to stress that the interaction term (\\ref{H_C_bosonized})\nis now quadratic in $\\phi_L$ and $\\phi_R$. Therefore the bosonization\napproach enables one to treat the charging energy exactly.\n\nFinally, the backscattering term (\\ref{backscattering_fermionic}) takes\nthe form\n\\begin{equation}\nH' = \\frac{D |r_L|}{\\pi} \\cos[2\\phi_L(0)]\n +\\frac{D |r_R|}{\\pi} \\cos[2\\phi_R(0)].\n \\label{backscattering_bosonized}\n\\end{equation}\nIn Eq.~(\\ref{backscattering_bosonized}) we neglected the products of\nMajorana fermions $\\eta_{L1}^{}\\eta_{L2}^{}$ and\n$\\eta_{R1}^{}\\eta_{R2}^{}$, as they commute with each other. Unlike the\nfirst two terms, $H_0$ and $H_C$, the backscattering term is not quadratic\nin bosonic variables and cannot be treated exactly. Below we assume weak\nbackscattering and develop the perturbation theory in small parameters\n$|r_L|$ and $|r_R|$.\n\n\n\\subsection{Perturbation theory for the conductance}\n\nIn this paper we are only interested in the linear conductance through the\nquantum dot. To find it, we will use a Kubo formula, which in the\nzero-frequency limit can be written as\n\\begin{equation}\nG = \\frac{i}{2\\hbar}\n \\int_{-\\infty}^{\\infty} t\n \\langle [ I(t), I(0) ] \\rangle dt,\n \\label{Kubo}\n\\end{equation}\nwhere $I$ is the operator of current through the dot.\n\nWe define $I$ as the average of the current $I_L$ flowing into the dot\nthrough the left contact and the current $I_R$ flowing out of the dot\nthrough the right contact. Clearly, the current $I_L$ is proportional to\nthe time derivative of the displacement of the effective elastic medium\ndescribing the left contact, $I_L\\sim e\\dot\\phi_L$, and, similarly,\n$I_R\\sim e\\dot\\phi_R$. From the equations of motion corresponding to the\nHamiltonian $H_0 + H_C + H'$, one concludes that $I\\sim\nev_F(\\Pi_L+\\Pi_R)$. A careful calculation gives\n\\begin{equation}\nI = \\frac{1}{2}ev_F\n \\left.\\left(\\Pi_L + \\Pi_R\\right)\\right|_{x=0}.\n \\label{current_operator}\n\\end{equation}\n\nAs the first step, we will find the conductance in the absence of the\nbarriers. Since the current operator (\\ref{current_operator}) commutes\nwith $H_C$, one should expect the conductance to be insensitive to the\ncharging energy. Clearly, at $E_C = 0$ the system is equivalent to two\nresistances $2\\pi\\hbar\/e^2$ of the two quantum point contacts, connected\nin series. One then expects the conductance to be\n\\begin{equation}\nG_0 = \\frac{e^2}{4\\pi\\hbar}.\n \\label{G_0}\n\\end{equation}\nThis result is easily derived from the Kubo formula (\\ref{Kubo}), because\nwith $H'=0$ the current-current correlator must be found for a quadratic\nHamiltonian. As a result, one finds $\\langle [ I(t), I(0) ] \\rangle =\ni(e^2\/2\\pi)\\delta'(t)$, and after the substitution into (\\ref{Kubo}) one\nreproduces (\\ref{G_0}).\n\nThe result (\\ref{G_0}) is independent of temperature. Since we used the\nHamiltonian (\\ref{H_0_fermionic}) with linearized spectra of electrons,\nthe temperature was assumed to be much smaller than the Fermi energy, but\ncould still be of the order of the charging energy. In the presence of the\nbarriers the conductance must acquire some temperature dependence, which\nis obvious from the limit of very high barriers (\\ref{cotunneling}).\nIndeed, the lowest non-vanishing (2nd-order) correction in $r_L$ and $r_R$\nto the conductance is already temperature dependent. At low temperatures\n$T\\ll E_C$, we get\n\\begin{equation}\nG=\\frac{e^2}{4\\pi\\hbar}\\!\\left(1-\\frac{\\pi\\Gamma_0(N)}{4T}\\right)\\!,\n \\label{2nd_order_spinless}\n\\end{equation}\nwhere $\\Gamma_0$ is defined as follows:\n\\begin{equation}\n\\Gamma_0(N)=\\frac{2\\gamma E_C}{\\pi^2}\n \\left[|r_L|^2+|r_R|^2+2|r_L||r_R|\\cos(2\\pi N)\\right]\\!.\n\\label{Gamma_0}\n\\end{equation}\nHere $\\gamma=e^{\\bf C}$, with ${\\bf C}\\approx 0.5772\\ldots$ being the\nEuler's constant. In Eq.~(\\ref{2nd_order_spinless}) we have neglected the\nterms small in $T\/E_C$. The result (\\ref{2nd_order_spinless}) and\n(\\ref{Gamma_0}) is analogous to that for the conductance of a 1D wire with\ntwo defects\\cite{Furusaki,remark3} and can be obtained in a similar way,\nsee Appendix~\\ref{A1}.\n\n\nAs expected, the presence of weak scattering in the contacts gives rise to\na small periodic correction to the conductance. The same behavior was\npredicted\\cite{Matveev1} for the average charge of the dot. However,\nunlike the average charge, the correction to conductance strongly depends\non temperature. As the temperature is lowered, the conductance decreases.\nThe only case when conductance is temperature-independent is when the\nbarriers are symmetric, $|r_L|=|r_R|$, and the dimensionless gate voltage\n$N$ is half-integer, which corresponds to the center of a peak in the\nhigh-barrier case. In this regime we expect the conductance to be of the\norder of $e^2\/\\hbar$, cf.\\ Sec.~\\ref{weak}, which is confirmed by the\nperturbation theory (\\ref{2nd_order_spinless}) and (\\ref{Gamma_0}).\n\n\n\\subsection{Non-perturbative calculation}\n\nAt low temperatures $T\\lesssim \\Gamma_0$ the second-order correction\nbecomes large, and the perturbation theory fails. To find the conductance\nin this case, one has to sum up infinite number of terms of the\nperturbation theory. To perform such a summation, we will use the\ntechnique developed in Ref.~\\onlinecite{Matveev1}.\n\nAt the first step we perform a linear transformation of the bosonic\nfields,\n\\begin{equation}\n\\phi_L = \\frac{\\phi_I - \\phi_C}{\\sqrt{2}},\n\\quad\n\\phi_R = \\frac{\\phi_I + \\phi_C}{\\sqrt{2}},\n \\label{transformation}\n\\end{equation}\nand similarly for $\\Pi_L$ and $\\Pi_R$. The advantage of using these new\nvariables is that the Coulomb energy (\\ref{H_C_bosonized}) is now\nexpressed in terms of the charging mode $\\phi_C$ only,\n\\begin{equation}\nH_C = \\frac{2E_C}{\\pi^2}\n \\left[\\phi_C(0)-\\frac{\\pi N}{\\sqrt{2}}\\right]^2.\n \\label{H_C_diagonalized}\n\\end{equation}\nThe charging energy (\\ref{H_C_diagonalized}) suppresses the fluctuations\nof $\\phi_C(0)$ at energy scales below $E_C$. Thus at $T\\ll E_C$ one can\nintegrate out the fluctuations of the charging mode by replacing the\nscattering term (\\ref{backscattering_bosonized}) with its value averaged\nover the fluctuations of $\\phi_C$. The resulting Hamiltonian is expressed\nin terms of the field $\\phi_I$ only and has the form\n\\begin{eqnarray}\nH_0 &=& \\frac{\\hbar v_F}{2}\\int\n \\left\\{\\frac{1}{\\pi}[\\nabla \\phi_I(x)]^2 + \\pi \\Pi_I^2(x)\n \\right\\}dx,\n \\label{H_0_simplified}\\\\\nH' &=& \\left(\\frac{\\gamma E_C D}{2\\pi^3}\\right)^{1\/2}\n \\left[re^{i\\sqrt{2}\\phi_I(0)}\n + r^{*}e^{-i\\sqrt{2}\\phi_I(0)}\\right],\n \\label{H'_simplified}\n\\end{eqnarray}\nwhere the dependence on the gate voltage is incorporated into the\ncomplex parameter $r = |r_L|e^{-i\\pi N} + |r_R|e^{i\\pi N}$.\n\nThe Hamiltonian (\\ref{H_0_simplified}) and (\\ref{H'_simplified}) can be\ntransformed to that of an impurity in a $g=\\frac12$ Luttinger\nliquid\\cite{Kane} by a simple transformation $\\phi=\\sqrt{2}\\phi_I$ and\n$\\Pi=\\Pi_I\/\\sqrt{2}$. The latter can be solved exactly\\cite{Kane} using\nthe technique developed in Ref.~\\onlinecite{Guinea}. For our purposes it\nwill be more convenient to use an alternative exact\nsolution.\\cite{Matveev1} The idea is to interpret $e^{i\\sqrt{2}\\phi_I}$ as\na bosonized form of some fermion annihilation operator.\\cite{Luther} This\ncan be done by rewriting the Hamiltonian in terms of two decoupled chiral\nbosonic fields defined as\n\\begin{equation}\n\\varphi_{\\pm}(x) = \\frac{\\phi_I(x) \\pm \\phi_I(-x)}{\\sqrt{2}}\n +\\pi\\int_{0}^{x}\n \\frac{\\Pi_I(x') \\pm \\Pi_I(-x')}{\\sqrt{2}} dx'.\n \\label{chiralization}\n\\end{equation}\nIn these variables the Hamiltonian takes the form\n\\begin{eqnarray}\nH_0 &=& \\frac{\\hbar v_F}{4\\pi}\\int\n \\left\\{[\\nabla \\varphi_{+}(x)]^2 + [\\nabla \\varphi_{-}(x)]^2\n \\right\\}dx,\n \\label{H_0_chiralized}\\\\\nH' &=& \\left(\\frac{\\gamma E_C D}{2\\pi^3}\\right)^{1\/2}\n \\left[re^{i\\varphi_{+}(0)}\n + r^{*}e^{-i\\varphi_{+}(0)}\\right].\n \\label{H'_chiralized}\n\\end{eqnarray}\nOne can easily check that the fields $\\varphi_+$ and $\\varphi_-$ commute.\nThus the $\\varphi_-$ part of the Hamiltonian completely decouples. Since\nthe field $\\varphi_+(x)$ satisfies the same commutation relations\n(\\ref{commutators}) as $\\varphi_{L1}$, one can in complete analogy with\nEq.~(\\ref{bosonization}) interpret $e^{i\\varphi_+}$ as a fermion\noperator,\\cite{remark}\n\n\\begin{figure}\n\\narrowtext\n\\epsfxsize=3in\n\\hspace{.2em}\n\\epsffile{4chan2.ps}\n\\vspace{1.5\\baselineskip}\n\\caption{Conductance (\\protect\\ref{conductance_spinless}) as a function\n of the dimensionless gate voltage $N$ for the symmetric case,\n $|r_L|=|r_R|=0.2$. The two curves are calculated for $E_C\/T=25$ (a) and\n $200$ (b). In this calculation we included the correction proportional\n to $T\/E_C$ in Eq.~(\\protect\\ref{finally_obtain}). The dashed line\n represents the conductance $G=G_LG_R\/(G_L + G_R)$ in the\n high-temperature limit, $T\\gg E_C$.}\n\\label{fig:2}\n\\end{figure}\n\\vspace{\\baselineskip}\n\n\\begin{equation}\ne^{i\\varphi_+(x)} = \\left(\\frac{2\\pi\\hbar v_F}{D}\\right)^{1\/2}\n \\eta_{+}\\psi_{+}(x).\n \\label{fermionization}\n\\end{equation}\n\nThe resulting Hamiltonian is quadratic in fermion operators,\n\\begin{eqnarray}\nH &=& i\\hbar v_F\\int\\psi_{+}^{\\dagger}(x)\\nabla\\psi_{+}(x)dx\n\\nonumber\\\\\n & & + \\left(\\frac{\\gamma \\hbar v_F E_C}{\\pi^2}\\right)^{1\/2}\n \\left[r\\eta_{+}^{}\\psi_{+}(0)\n + r^{*}\\psi_{+}^{\\dagger}(0)\\eta_{+}^{}\n \\right],\n \\label{H_fermionized}\n\\end{eqnarray}\nand can be easily diagonalized.\\cite{Matveev1} This greatly simplifies the\ncalculation of the conductance. One can again use the Kubo formula\n(\\ref{Kubo}), with the properly transformed current operator\n(\\ref{current_operator}). In the new fermionic variables the current\noperator has the form\n\\begin{equation}\nI = -e v_F \\left. :\\!\\psi_{+}^{\\dagger}\\psi_{+}\\!\\!:\n \\right|_{x=0}.\n \\label{current_operator_fermionized}\n\\end{equation}\n\nThe actual calculation of the linear conductance is rather long, but\nstraightforward. First we find the current-current correlation function\nto be\n\\begin{eqnarray}\n& &\\langle [I(t),I(0)]\\rangle =\n i\\frac{e^2}{2\\pi}\\delta'(t)\n\\nonumber\\\\\n & &\\hspace{3em}\n - i\\frac{e^2}{2\\pi^2\\hbar^2}\\int\\!\\!\\!\\int\\!\n \\frac{\\Gamma_0^2\\sin[(\\epsilon+\\epsilon')t\/\\hbar]d\\epsilon d\\epsilon'}\n {(\\epsilon^2+\\Gamma_0^2)(e^{\\epsilon\/T}+1)(e^{\\epsilon'\/T}+1)}.\n \\label{current-current}\n\\end{eqnarray}\nThe substitution of Eq.~(\\ref{current-current}) into the Kubo formula\n(\\ref{Kubo}) gives the following expression for the linear conductance:\n\\begin{equation}\nG=\\frac{e^2}{4\\pi\\hbar}\n \\left[1-\\int^\\infty_{-\\infty}dE\n \\frac{1}{4T\\cosh^2\\frac{E}{2T}}\n \\frac{\\Gamma^2_0(N)}{E^2+\\Gamma^2_0(N)}\n \\right].\n \\label{conductance_spinless}\n\\end{equation}\nEquation (\\ref{conductance_spinless}) is the central result of this\nsection. It is valid for any temperature below $E_C$, and, unlike the\nperturbative result (\\ref{2nd_order_spinless}), it is not restricted to\nonly sufficiently high temperatures $T\\gg\\Gamma_0$. Of course, at\n$T\\gg\\Gamma_0$ the result (\\ref{2nd_order_spinless}) is reproduced.\n\nIn the opposite limit, $T\\ll\\Gamma_0$, the conductance is quadratic in\ntemperature,\n\\begin{equation}\nG = \\frac{\\pi e^2}{12\\hbar}\\left[\\frac{T}{\\Gamma_0(N)}\\right]^2.\n \\label{low_T_spinless}\n\\end{equation}\nIt is important to note that the same temperature dependence was found for\nthe inelastic co-tunneling (\\ref{cotunneling}) in the weak-tunneling\nregime. As expected, the quadratic temperature dependence is a universal\nproperty of inelastic co-tunneling and is not limited to the\nweak-tunneling limit, cf.\\ Eq.~(\\ref{tails}).\n\nIn the special case of symmetric barriers, $|r_L| = |r_R| = r_0$, the\nlow-temperature limit (\\ref{low_T_spinless}) is never achieved on\nresonance. Indeed, we have $\\Gamma_0= (8\\gamma E_C\/\\pi^2)r_0^2\\cos^2 \\pi\nN$, i.e., $\\Gamma_0$ vanishes at half-integer values of $N$. Therefore in\nthe low-temperature limit the resonant value of conductance is\n$e^2\/4\\pi\\hbar$, whereas away from the resonance $\\Gamma_0\\neq 0$, and\n$G\\propto T^2$. It is clear then that at low temperatures the weak Coulomb\nblockade oscillation in conductance (\\ref{2nd_order_spinless}) transforms\ninto sharp periodic peaks, Fig.~\\ref{fig:2}. Unlike the weak-tunneling\ncase (\\ref{cotunneling}), the tails of the peak at $N=\\frac12$ are not\nLorentzian: $G\\propto1\/(N-\\frac12)^4$.\n\nFinally, it is worth noting that the conductance\n(\\ref{conductance_spinless}) coincides with that of a weak impurity in the\n$g=\\frac12$ Luttinger liquid.\\cite{Kane} This is not surprising,\nconsidering the aforementioned identity of the Hamiltonian\n(\\ref{H_0_simplified}), (\\ref{H'_simplified}) with the one for the quantum\nimpurity problem.\\cite{Kane} In the latter problem the condition\n$g=\\frac12$ requires special fine-tuning of the strength of interaction\nbetween electrons. In our case the condition $g=\\frac12$ is satisfied\nautomatically when we integrate out one of the two bosonic\nmodes.\\cite{Furusaki}\n\n\n\\section{Co-tunneling of electrons with spins}\n\\label{withspin}\n\nIn the previous section we considered the case of spinless electrons,\nwhich can be realized by applying strong in-plane magnetic field. In the\nabsence of the field, each of the contacts has two identical modes,\ncorresponding to two possible values of electron spin. The number of modes\nis an important parameter in our problem. This can be easily seen from the\nanalogy to the Kondo problem discussed in Sec.~\\ref{weak}. In the\nweak-tunneling case, the model of spinless electrons maps to the 2-channel\nKondo problem, because the fictitious spin of the impurity---i.e., the\ncharge of the dot---can be changed by tunneling of an electron in two\nindependent channels, corresponding to the left and right contacts. In the\npresence of real spins of electrons, the number of channels in the\neffective Kondo problem is 4. It is well known that the number of channels\nin the Kondo problem strongly affects its properties.\\cite{NB} Thus one\nshould expect that the presence of spins should have a dramatic effect\nupon the transport through a quantum dot.\n\n\n\\subsection{Strong-tunneling case}\n\nWe will start with the discussion of the strong-tunneling limit, ${\\cal\n T}_L, {\\cal T}_R \\to 1$. To find the conductance, one can use the same\ntechnique as in Sec.~\\ref{spinless}. The Hamiltonian of the problem can be\nwritten in the bosonized form, analogous to\n(\\ref{H_0_bosonized})--(\\ref{backscattering_bosonized}),\n\\begin{eqnarray}\nH_0 &=& \\frac{\\hbar v_F}{2}\\!\\sum_{\\alpha}\\int\\!\n \\left\\{\\frac{1}{\\pi}[\\nabla \\phi_{L\\alpha}(x)]^2\n + \\pi \\Pi_{L\\alpha}^2(x)\\right.\n\\nonumber\\\\\n &&\\hspace{5.3em}\\left. +\\frac{1}{\\pi}[\\nabla \\phi_{R\\alpha}(x)]^2\n + \\pi \\Pi_{R\\alpha}^2(x)\n \\right\\}\\!dx,\n\\label{H_0_with_spins}\\\\\nH_{C} &=& \\frac{E_C}{\\pi^2}\\left\\{\\sum_{\\alpha}\n [\\phi_{R\\alpha}(0) - \\phi_{L\\alpha}(0)] -\\pi N\\right\\}^2,\n\\label{H_C_with_spins}\\\\\nH' &=& \\frac{D}{\\pi} \\sum_{\\alpha}\n \\{|r_L| \\cos[2\\phi_{L\\alpha}(0)]\\!\n +\\! |r_R| \\cos[2\\phi_{R\\alpha}(0)]\n \\}.\n \\label{backscattering_with_spins}\n\\end{eqnarray}\nHere the fields with $\\alpha=\\ \\uparrow$ and $\\downarrow$ describe the\nelectrons with spins up and down respectively.\n\nOne can again find the linear conductance using the Kubo formula\n(\\ref{Kubo}), but the current operator (\\ref{current_operator}) must be\nmodified to account for the spins,\n\\begin{equation}\nI = \\frac{1}{2}ev_F\\!\\sum_{\\alpha=\\uparrow,\\downarrow}\\!\n \\left(\\Pi_{L\\alpha} + \\Pi_{R\\alpha}\\right)\\Bigl|_{x=0}.\n \\label{current_operator_with_spins}\n\\end{equation}\n\nIn the absence of barriers, $|r_L|=|r_R|=0$, the calculation is trivial,\nas the Hamiltonian (\\ref{H_0_with_spins}) and (\\ref{H_C_with_spins})\nis quadratic. Due to the doubling of the number of modes in each contact,\nin the presence of spins the conductance doubles compared to\nEq.~(\\ref{G_0}),\n\\begin{equation}\nG_0 = \\frac{e^2}{2\\pi\\hbar}.\n \\label{G_0_with_spins}\n\\end{equation}\n\nWeak backscattering, $|r_L|,|r_R|\\ll1$, can be treated in perturbation\ntheory in $H'$, see Appendix \\ref{A2}. Up to the second order we get\n\\begin{equation}\n G = \\frac{e^2}{2\\pi\\hbar}\n \\left[1-\\frac{2\\Gamma(\\frac{3}{4})}{\\Gamma(\\frac{1}{4})}\n \\sqrt{\\frac{\\gamma E_C}{\\pi T}}\n \\left(|r_L|^2+|r_R|^2\\right)\n \\right].\n \\label{2nd_order_with_spins}\n\\end{equation}\nIt is instructive to compare this result with its analog for the spinless\ncase, namely Eq.~(\\ref{2nd_order_spinless}). In both cases the small\ncorrection to the conductance has the tendency to grow as the temperature\nis decreased. An important difference is that the second-order correction\nin (\\ref{2nd_order_with_spins}) does not depend on the gate voltage $N$,\nand as a result it grows even at half-integer $N$, i.e., on resonance.\nThus, in contrast to the spinless case, one should expect the resonance\nvalue to be smaller than the unperturbed conductance $G_0=e^2\/2\\pi\\hbar$.\n\nThis difference has a natural interpretation in terms of the analogy with\nthe Kondo model. As we already mentioned, the spinless case corresponds to\nthe 2-channel Kondo model. In the 2-channel case the strong-coupling fixed\npoint is stable at zero magnetic field, or, in terms of the Coulomb\nblockade problem, at half-integer $N$. Therefore, precisely on resonance a\nweak reflection must be an irrelevant perturbation, and the peak value of\nthe conductance should be given by Eq.~(\\ref{G_0}). On the contrary, the\n4-channel Kondo problem corresponding to the case with spins has a stable\nfixed point at an intermediate value of coupling, and both weak- and\nstrong-coupling fixed points are unstable.\\cite{NB} Thus even on resonance\nthe conductance is not given by its value in the strong-tunneling limit\n(\\ref{G_0_with_spins}), as confirmed by the perturbation theory\n(\\ref{2nd_order_with_spins}).\n\n\\subsection{Conductance in the case of asymmetric barriers}\n\\label{asymmetric}\n\nUnfortunately, it is not easy to calculate the conductance near the\nintermediate-coupling fixed point. Some conclusions about it can be drawn\non the basis of the scaling approach, Sec.~\\ref{scaling}. One should note,\nhowever, that this fixed point is stable only if the coupling in all the 4\nchannels is identical. In the case of the Coulomb blockade problem, the\ncoupling in the pairs of channels corresponding to different spin\norientations in the same contact must be identical (in the absence of\nmagnetic field). On the other hand, the coupling constants for the modes\nin different contacts are determined by the corresponding conductances and\nare not necessarily equal. In fact, in the structure shown in\nFig.~\\ref{fig:1}, the two conductances are controlled separately by the\ngate voltages $V_L$ and $V_R$. Therefore in experiments the coupling\nconstants corresponding to the left and right modes are not equal, unless\na special attempt to make them equal is made.\n\nIn this section we study the limit of very asymmetric barriers, namely we\nassume that one of the barriers is very high, whereas the other one is\nlow. In terms of the transmission coefficients of the barriers this limit\ncorresponds to ${\\cal T}_L \\ll 1$ and $1-{\\cal T}_R \\ll 1$. Such a regime\ncan be easily realized in an experiment by proper tuning of the gate\nvoltages $V_L$ and $V_R$. In addition, we will see in Sec.~\\ref{scaling}\nthat at $T\\to 0$ arbitrarily weak asymmetry grows and the system scales\ntowards the strongly asymmetric limit. Therefore at low temperatures the\nresults of this section can be applied to any asymmetric system.\n\nIn Sec.~\\ref{weak} we described the contact in the regime of weak\ntunneling by the tunnel Hamiltonian (\\ref{tunnel_hamiltonian}) written in\nterms of the original electrons. On the other hand, in Sec.~\\ref{spinless}\nit was more convenient to use the bosonized form of the Hamiltonian for\nthe electron gas in the contact with transmission coefficient close to\nunity. Therefore when considering the asymmetric case ${\\cal T}_L\\to0$ and\n${\\cal T}_R\\to1$ we shall bosonize only the electrons in the right\ncontact. Then the unperturbed fixed-point Hamiltonian $H_0 + H_C$ is given\nby\n\\begin{eqnarray}\nH_0 &=& \\sum_{k\\alpha} \\epsilon_k a_{k\\alpha}^\\dagger a_{k\\alpha}^{}\n + \\sum_{p\\alpha} \\epsilon_p a_{p\\alpha}^\\dagger a_{p\\alpha}^{}\n\\nonumber\\\\\n & & +\\frac{\\hbar v_F}{2}\\!\\sum_{\\alpha}\\!\\int\\!\n \\left\\{\\frac{1}{\\pi}[\\nabla \\phi_{R\\alpha}(x)]^2\n + \\pi \\Pi_{R\\alpha}^2(x)\n \\right\\}dx,\n \\label{H_0_asymmetric}\\\\\nH_C &=& E_C\\!\n \\left[\\hat n+\\frac{1}{\\pi}\\sum_\\alpha\\phi_{R\\alpha}(0)\n - N\\right]^2.\n \\label{H_C_asymmetric}\n\\end{eqnarray}\nHere the operators $a_{k\\alpha}$ and $a_{p\\alpha}$ correspond to the\nleft-contact electrons in the lead and in the dot, respectively; the\nbosonic fields $\\phi_{R\\alpha}$ and $\\Pi_{R\\alpha}$ model the electrons in\nthe right contact. An integer-valued operator $\\hat n$ describes the number\nof electrons transferred to the dot through the left barrier.\n\nThe perturbation consists of two parts: weak tunneling through the high\nbarrier in the left contact, and weak scattering on the small barrier in\nthe right contact,\n\\begin{eqnarray}\n H_L' &=& \\sum_{kp\\alpha}\n (v_t^{} a_{k\\alpha}^\\dagger a_{p\\alpha}^{}F\n + v_t^* a_{p\\alpha}^\\dagger a_{k\\alpha}^{}F^\\dagger),\n\\label{H_L'}\\\\\n H_R' &=& \\frac{1}{\\pi}D|r_R| \\sum_{\\alpha}\n \\cos[2\\phi_{R\\alpha}(0)].\n\\label{H_R'}\n\\end{eqnarray}\nWe have chosen not to write the operator $\\hat n$ of the number of\nelectrons transferred into the dot through the barrier in terms of the\nfermionic operators as $\\hat n=\\sum a_{p\\alpha}^\\dagger a_{p\\alpha}^{}$,\nbut to treat it as an independent variable. Thus the tunneling Hamiltonian\n(\\ref{H_L'}) acquires the charge-lowering and raising operators $F$ and\n$F^\\dagger$ defined as\n\\begin{equation}\n[F,\\hat n] = F.\n \\label{lowering_operator}\n\\end{equation}\nOne possible representation\\cite{Ingold} of these operators is given by\n$F=e^{ix}$ and $\\hat n = i\\frac{\\partial}{\\partial x}$, where $x$ has the\nmeaning of a phase operator conjugated to $\\hat n$.\n\nOur goal in this section is to find the conductance of the system,\nassuming that the perturbations $H_L'$ and $H_R'$ are small. At the first\nstep we will utilize the smallness of the amplitude $v_t^{}$ of tunneling\nthrough the left barrier, and find the expression for the current in the\nsecond order in $v_t^{}$. The current can be defined as the time\nderivative of the number of electrons that have tunneled through the\nbarrier, $I = -e \\frac{\\partial \\hat n}{\\partial t} =\n\\frac{ie}{\\hbar}[\\hat n,H]$. Using the commutation relation\n(\\ref{lowering_operator}), we get\n\\begin{equation}\nI = i\\frac{e}{\\hbar}\n \\sum_{kp\\alpha} \\left(\n v_t^* a_{p\\alpha}^\\dagger a_{k\\alpha}^{} F^\\dagger -\n v_t a_{k\\alpha}^\\dagger a_{p\\alpha}^{} F\\right).\n \\label{current_asymmetric}\n\\end{equation}\nTo find the conductance through the dot, we will substitute the above\ncurrent operator into the Kubo formula (\\ref{Kubo}). This requires the\ncalculation of the correlator $\\langle[I(t),I(0)]\\rangle$. Since $I\\propto\n|v_t|$, up to the second order in tunneling matrix elements we can average\n$[I(t),I(0)]$ over the equilibrium thermal distribution of the system\nwithout tunneling, $H_L'=0$. Then the averaging over the fermionic degrees\nof freedom is straightforward, and we get\n\\begin{equation}\nG = -i G_L \\frac{\\pi T^2}{\\hbar^2}\n \\int_{-\\infty}^{\\infty}\n \\frac{tK(t)dt}{\\sinh^2[\\pi T(t-i\\delta)\/\\hbar]}.\n \\label{conductance_asymmetric}\n\\end{equation}\nHere we have introduced the conductance of the left tunnel barrier,\n$G_L=(4\\pi e^2\/\\hbar)\n\\sum_{kp}|v_t^{}|^2\\delta(\\epsilon_k)\\delta(\\epsilon_p)$, and the\ncorrelator\n\\begin{equation}\nK(t) = \\langle F(t) F^\\dagger(0)\\rangle\n = \\langle F^\\dagger(t) F(0)\\rangle.\n \\label{K}\n\\end{equation}\nThe equality of the two correlators in Eq.~(\\ref{K}) follows from the fact\nthat the symmetry transformation $\\hat n \\to -\\hat n$, $\\phi_{R\\alpha}\\to\n-\\phi_{R\\alpha}$, and $N\\to -N$ does not affect the Hamiltonian\n(\\ref{H_0_asymmetric}), (\\ref{H_C_asymmetric}), and (\\ref{H_R'}), but\nchanges $F^\\dagger\\leftrightarrow F$.\n\nIt is instructive to apply the formula (\\ref{conductance_asymmetric}) to\nthe non-interacting case, $E_C = 0$. In this limit the operators $F$ and\n$F^\\dagger$ commute with the Hamiltonian (\\ref{H_0_asymmetric}) and\n(\\ref{H_R'}). Thus the correlator (\\ref{K}) is unity, and we get $G=G_L$.\nThis is the expected result, since we assumed $G_L\\to0$, and the\nconductance of the system must be determined by that of the weakest link.\n\nLet us now find the conductance at non-zero charging energy $E_C$, but in\nthe absence of the scattering in the right contact, $r_R=0$. In this\ncase the correlator (\\ref{K}) is no longer trivial, because the operators\n$F$ and $F^\\dagger$ do not commute with the interaction term\n(\\ref{H_C_asymmetric}). When an electron tunnels into the dot, $n\\to n+1$,\nthe system is shifted from its equilibrium state. To accommodate the\nadditional electron brought by the tunneling process and to return the\nsystem to the equilibrium state, one electron must leave the dot\nthrough the right contact. Thus the time evolution of the operators $F$\nand $F^\\dagger$ becomes non-trivial.\n\nTo find the correlator $K(t)$, we will perform the following unitary\ntransformation of the Hamiltonian:\n\\begin{eqnarray}\n&U = e^{i\\hat n \\Theta},&\n \\label{unitary_transformation}\\\\\n&{\\displaystyle\n\\Theta = \\frac{\\pi}{2} \\int_{-\\infty}^{\\infty}\n \\left[\\Pi_{R\\uparrow}(y) + \\Pi_{R\\downarrow}(y)\\right]dy}.&\n \\label{phase}\n\\end{eqnarray}\nThis transformation shifts the bosonic fields $\\phi_{R\\alpha}$ and adds a\nphase factor to the charge-lowering operator $F$,\n\\begin{eqnarray}\n \\phi_{R\\alpha}(x)&\\to& \\phi_{R\\alpha}(x) - \\frac{\\pi}{2}\\hat n,\n\\label{transformed_phi}\\\\\n F&\\to& F e^{i\\Theta}.\n\\label{transformed_F}\n\\end{eqnarray}\nOne can easily see that the kinetic-energy term (\\ref{H_0_asymmetric}) is\ninvariant under the transformation (\\ref{unitary_transformation}), whereas\nthe charging term (\\ref{H_C_asymmetric}) transforms to a $\\hat\nn$-independent form,\n\\begin{equation}\n\\tilde H_C = U^\\dagger H_C U =\nE_C\\left[\\frac{1}{\\pi}\\sum_\\alpha\\phi_{R\\alpha}(0) - N\\right]^2.\n\\label{transformed_H_C}\n\\end{equation}\nUpon this transformation, the operators $F$ and $F^\\dagger$ commute with\nthe Hamiltonian (\\ref{H_0_asymmetric}) and (\\ref{transformed_H_C}), but\nthey acquire non-trivial phase factors, see Eq.~(\\ref{transformed_F}).\nThis enables us to find $K(t)$ as the correlator of the phase factors,\n$K(t) = K_{\\Theta}(t)\\equiv \\langle e^{i\\Theta(t)}\ne^{-i\\Theta(0)}\\rangle$. Since the phase $\\Theta$ is linear in bosonic\nvariables, the calculation of the correlator $K_{\\Theta}$ is\nstraightforward:\n\\begin{eqnarray}\nK_{\\Theta}(t)\n &=& \\exp\\{-\\langle[\\Theta(0)-\\Theta(t)]\\Theta(0)\\rangle\\}\n\\nonumber\\\\\n &\\simeq& \\frac{\\pi^2 T}{2i\\gamma E_C}\n \\frac{1}{\\sinh[\\pi T(t-i\\delta)\/\\hbar]}.\n \\label{Debye_Waller}\n\\end{eqnarray}\nThe last equality is valid asymptotically at low energies, $T,\n\\hbar\/t\\ll E_C$. The substitution of $K=K_\\Theta$ into the\nexpression for the conductance (\\ref{conductance_asymmetric}) leads to a\nlinear temperature dependence of the conductance:\n\\begin{equation}\nG = G_L \\frac{\\pi^3 T}{8\\gamma E_C}.\n \\label{peak_conductance_asymmetric}\n\\end{equation}\n\nPhysically the origin of the linear suppression of the conductance at low\ntemperatures (\\ref{peak_conductance_asymmetric}) is the orthogonality\ncatastrophe.\\cite{Anderson2} Unlike the non-interacting case $E_C = 0$\ndiscussed above, the charge fluctuations in the dot are now suppressed at\nenergies below $E_C$. Effectively one can interpret the effect of the\ncharging energy as a hard-wall boundary condition for the wavefunctions of\nthe right-contact electrons with energies within the band of width $W\\sim\nE_C$ near the Fermi level. When an electron tunnels through the left\nbarrier, the system must lower its charging energy by moving one electron\nthrough the right contact. Since the two spin channels are completely\nsymmetric, each mode transfers the charge $q=\\frac{e}{2}$. According to\nthe Friedel sum rule, this leads to a change in the boundary condition\ncorresponding to an additional scattering phase shift $\\delta=\\pm\\frac{\\pi\n q}{e}=\\pm\\frac{\\pi}{2}$. A sudden change of the boundary conditions\ncreates a state with a large number of electron-hole excitations, which is\nnearly orthogonal\\cite{Anderson2} to the ground state of the system. This\northogonality leads to a power-law suppression of the tunneling density of\nstates $\\nu(\\varepsilon)\\propto\\varepsilon^\\chi$, with the\nexponent\\cite{Schotte,ND,Mahan} determined by the phase shifts in all the\nelectronic channels, $\\chi=\\sum(\\delta\/\\pi)^2$. In our system, each spin\nmode gives two channels---in the dot and in the right lead---leading to\nthe total of 4 channels. Thus the exponent $\\chi$ is unity, and the\ndensity of states $\\nu(\\varepsilon)\\propto\\varepsilon$. The linear\nsuppression of the tunneling density of states results in the $T$-linear\nconductance (\\ref{peak_conductance_asymmetric}).\n\nThe linear temperature dependence of conductance\n(\\ref{peak_conductance_asymmetric}) appears to contradict to the $T^2$ law\nfor the inelastic co-tunneling, Eqs.~(\\ref{cotunneling}), (\\ref{tails})\nand (\\ref{low_T_spinless}). The reason is the absence of the barrier in\nthe right junction assumed in the derivation of\nEq.~(\\ref{peak_conductance_asymmetric}). We will now show that the\npresence of arbitrarily weak barrier in the right contact is crucial (see\nalso Appendix~\\ref{A3}) and leads to the quadratic temperature dependence\nof the conductance at $T\\to0$.\n\nThe weak backscattering in the right contact is described by the term\n(\\ref{H_R'}) in the Hamiltonian. As it follows from\nEq.~(\\ref{transformed_phi}), under the unitary transformation\n(\\ref{unitary_transformation}) the scattering term (\\ref{H_R'}) acquires\nadditional sign:\n\\begin{equation}\n\\tilde H_R' = (-1)^{\\hat n}\n \\frac{1}{\\pi}D|r_R| \\sum_{\\alpha}\n \\cos[2\\phi_{R\\alpha}(0)].\n \\label{transformed_H_R'}\n\\end{equation}\nThe operators $F$ and $F^\\dagger$ no longer commute with the Hamiltonian:\nthey commute with $H_0 + \\tilde H_C$ and anti-commute with $\\tilde H_R'$.\nTherefore the correlator (\\ref{K}) is no longer equal to $K_{\\Theta}$, but\nrather $K(t)=K_{\\Theta}(t)K_{F}(t)$, where\n\\begin{equation}\nK_F(t) = \\langle F(t) F^\\dagger(0)\\rangle\n \\label{K_F}.\n\\end{equation}\nHere the averaging $\\langle\\ldots\\rangle$ is over the equilibrium thermal\ndistribution of the transformed Hamiltonian, $H=H_0+\\tilde H_C + \\tilde\nH_R'$.\n\nIn the absence of the backscattering in the right contact, $K_{\\Theta}$ is\ngiven by Eq.~(\\ref{Debye_Waller}), and $K_F(t)\\equiv 1$. When a weak\nbackscattering is added, both correlators are affected. It is clear,\nhowever, that $K_\\Theta$ is not modified significantly. Indeed, it follows\nfrom the above mentioned analogy with the orthogonality catastrophe that\nthe time dependence (\\ref{Debye_Waller}) is determined solely by the\ncharge transferred through the right junction, which in our case is always\n$e$---the charge of the electron brought into the dot through the left\nbarrier---and is not affected by the presence of the right barrier.\nTherefore the only change in $K_\\Theta(t)$ caused by the weak barrier can\nbe a small modification of the prefactor in Eq.~(\\ref{Debye_Waller}). On\nthe other hand, we will show that $K_F(t)$ is modified significantly,\nwhich leads to the quadratic temperature dependence of the conductance.\n\nTo find $K_F(t)$ we first note that the Hamiltonian $H=H_0+\\tilde H_C +\n\\tilde H_R'$ is identical to the one considered in Sec.~\\ref{spinless} and\nin Ref.~\\onlinecite{Matveev1}. The difference with the problem considered\nin Sec.~\\ref{spinless} is that instead of two modes corresponding to two\ndifferent contacts we now have two spin modes in the same (right) contact.\nThus the two scattering amplitudes are equal and given by $|r_R|$.\nAs we have seen in Sec.~\\ref{spinless}, the presence of the backscattering\nshows up at low-energy scales $\\epsilon\\sim |r_R|^2 E_C\\ll E_C$, where one\ncan effectively integrate out the mode related to the charge fluctuations\nin the dot, $\\phi_{R\\uparrow} + \\phi_{R\\downarrow}$. The resulting\nHamiltonian is expressed in terms of the spin field\n$\\phi_{Rs}=(\\phi_{R\\uparrow} - \\phi_{R\\downarrow})\/\\sqrt{2}$ and has the\nform:\n\\begin{eqnarray}\nH &=& \\frac{\\hbar v_F}{2}\\!\\int\\!\n \\left\\{\\frac{1}{\\pi}[\\nabla \\phi_{Rs}(x)]^2\n + \\pi \\Pi_{Rs}^2(x)\n \\right\\}dx\n\\nonumber\\\\\n & & + (-1)^{\\hat n}\n \\sqrt{\\frac{8\\gamma E_C D}{\\pi^3}}\\,|r_R|\\cos\\pi N\n \\cos\\!\\left[\\sqrt{2}\\phi_{Rs}(0)\\right]\\!.\n \\label{spin_Hamiltonian}\n\\end{eqnarray}\nWe can now fermionize the Hamiltonian using a representation of the\noperators $e^{\\pm i\\sqrt{2}\\phi_{Rs}}$ in terms of the fermion creation\nand annihilation operators in a way identical to the one leading from the\nHamiltonian (\\ref{H_0_simplified}) and (\\ref{H'_simplified}) to\nEq.~(\\ref{H_fermionized}). The final Hamiltonian is again quadratic in\nfermion operators, but also depends on the operator $\\hat n$ of the number\nof particles transferred through the left barrier,\n\\begin{eqnarray}\nH &=& i\\hbar v_F\\int\\psi^{\\dagger}(x)\\nabla\\psi(x)dx\n\\nonumber\\\\\n & & + (-1)^{\\hat n}\n \\frac{2}{\\pi}\\sqrt{\\gamma \\hbar v_F E_C}\\,\n |r_R|\\cos\\pi N\n\\nonumber\\\\\n & &\\hspace{1em}\\times\n (c+c^\\dagger)[\\psi(0) - \\psi^{\\dagger}(0)].\n \\label{H_fermionized_asymmetric}\n\\end{eqnarray}\nHere we presented a Majorana fermion operator similar to $\\eta_+$ in\nEq.~(\\ref{H_fermionized}) as a sum of fermion creation and annihilation\noperators, $\\eta=c+c^\\dagger$. This representation is more convenient,\nbecause we can now replace the operators $F$ and $F^\\dagger$ in the\ndefinition (\\ref{K_F}) of the correlator $K_F$ by\n$F_s=F_s^\\dagger=1-2c^\\dagger c$. Indeed, these new operators commute with\nthe first term in the Hamiltonian (\\ref{H_fermionized_asymmetric}),\nanti-commute with the second one, and possess the property\n$F_s^{}F_s^\\dagger = FF^\\dagger=1$. Therefore one can find the correlator\n$K_F$ as\\cite{remark2}\n\\begin{equation}\nK_F(t) = \\langle[1-2c^\\dagger(t)c(t)]\n [1-2c^\\dagger(0)c(0)]\\rangle.\n \\label{K_F_fermionized}\n\\end{equation}\n\nThe actual calculation of the correlator (\\ref{K_F_fermionized}) is now\nstraightforward, as the Hamiltonian (\\ref{H_fermionized_asymmetric}) is\nquadratic in the fermion operators. The result for the correlator $K_F$\nhas the form\n\\begin{equation}\nK_F(t) = \\frac{2\\Gamma_R}{\\pi}\\int_{-\\infty}^{\\infty}\n \\frac{e^{iEt}dE}\n {(E^2+\\Gamma_R^2)(e^{E\/T}+1)}.\n \\label{K_F_final}\n\\end{equation}\nHere we have introduced the low energy scale\n\\begin{equation}\n\\Gamma_R = \\frac{8\\gamma}{\\pi^2}E_C |r_R|^2 \\cos^2 \\pi N,\n \\label{Gamma_R}\n\\end{equation}\noriginating from the presence of a weak scatterer in the right contact. In\nthe weak-scattering limit, $\\Gamma_R\\to0$, Eq.~(\\ref{K_F_final})\nreproduces the result $K_F(t)=1$ mentioned above. However, in\nthe opposite limit $T,\\hbar\/t\\ll\\Gamma_R$, the correlator $K_F$ has the\nsame non-trivial time dependence as the correlator $K_{\\Theta}(t)$\ngiven by Eq.~(\\ref{Debye_Waller}),\n\\[\nK_F(t) \\simeq \\frac{2T}{i\\Gamma_R}\n \\frac{1}{\\sinh[\\pi T (t -i\\delta)\/\\hbar]}.\n\\]\nThis additional time-dependent contribution to $K(t)$ leads to the\nquadratic temperature dependence of the linear conductance at\n$T\\ll\\Gamma_R$.\n\nTo find the conductance at temperatures $T\\sim\\Gamma_R$, we can now\nsubstitute the correlator $K(t)=K_\\Theta(t)K_F(t)$ into\nEq.~(\\ref{conductance_asymmetric}), which leads to the following\nexpression:\n\\begin{equation}\nG = \\frac{TG_L}{8\\gamma E_C}\n \\int^\\infty_{-\\infty}\n \\frac{\\Gamma_R}{E^2+\\Gamma^2_R}\n \\frac{\\pi^2+(E\/T)^2}{\\cosh^2(E\/2T)}dE.\n \\label{conductance_asymmetric_final}\n\\end{equation}\nEquation (\\ref{conductance_asymmetric_final}) is the central result of\nthis section. It is valid for any $T\/\\Gamma_R$, if both temperature and\n$\\Gamma_R$ are much smaller than the charging energy $E_C$. The condition\n$\\Gamma_R\\ll E_C$ means $|r_R|\\ll1$, i.e., the scattering in the right\ncontact must indeed be weak.\n\n\n\n\\begin{figure}\n\\narrowtext\n\\epsfxsize=2.8in\n\\hspace{.4em}\n\\epsffile{4chan3.ps}\n\\vspace{1.5\\baselineskip}\n\\caption{Conductance (\\protect\\ref{conductance_asymmetric_final})\n as a function of the dimensionless gate voltage $N$ for the asymmetric\n case, $|r_R|=0.4$. The three curves are calculated for\n $E_C\/T=10$ (a), $40$ (b), and $160$ (c).}\n\\label{fig:3}\n\\end{figure}\n\\vspace{\\baselineskip}\n\nIn the limit $T\\gg\\Gamma_R$, Eq.~(\\ref{conductance_asymmetric_final})\nreproduces the previously obtained linear temperature dependence\n(\\ref{peak_conductance_asymmetric}) in the absence of the scattering in\nthe right contact. It is important to note, however, that the linear\ntemperature dependence is obtained not only at $r_R=0$, but also in the\npresence of the barrier if $N$ is half-integer, i.e., at the centers of\nconductance peaks.\n\nAway from the centers of the peaks, the conductance has a stronger\ntemperature dependence, $G\\propto T^2$. Indeed, at $T\\ll\\Gamma_R$ the\nexpression (\\ref{conductance_asymmetric_final}) reduces to\n\\begin{equation}\nG = \\frac{2\\pi^2}{3\\gamma}\\frac{T^2}{E_C\\Gamma_R} G_L.\n \\label{off-peak_conductance}\n\\end{equation}\nThus the quadratic temperature dependence of the linear conductance is\nrestored for the off-peak values of the gate voltage. The evolution of the\nconductance peaks is illustrated in Fig.~\\ref{fig:3}.\n\n\n\n\\section{Scaling approach to inelastic co-tunneling}\n\\label{scaling}\n\nIn the sections \\ref{weak}--\\ref{withspin} we have considered several\ncases when the inelastic co-tunneling can be studied analytically even\nthough one or both contacts are in the strong-tunneling regime. A common\nproperty of all these results is that at $T\\to0$ the conductance behaves\nas $T^2$ away from the resonance values of the gate voltage, i.e., at\n$N\\neq\\pm\\frac12, \\pm\\frac32,\\ldots$. The same temperature dependence was\nobtained earlier\\cite{AN} for the off-resonance conductance in the\nweak-tunneling case, see Eq.~(\\ref{cotunneling}). In this section we study\nthe general scaling properties of our system and argue that the quadratic\ntemperature dependence of linear conductance is universal and should hold\nfor arbitrary barriers.\n\nWe will describe our system by the bosonized Hamiltonian\n(\\ref{H_0_bosonized})--(\\ref{backscattering_bosonized}) in the spinless\ncase and (\\ref{H_0_with_spins})--(\\ref{backscattering_with_spins}) in the\npresence of spins. In the absence of the backscattering the Hamiltonian is\nquadratic, i.e., $|r_L|=|r_R|=0$ is a fixed point of our Hamiltonian. In\nanalogy with the Kondo problem we will call it the strong-coupling fixed\npoint. To investigate the stability of this fixed point, one should find\nthe scaling dimension of the backscattering term\n(\\ref{backscattering_bosonized}) or (\\ref{backscattering_with_spins}).\n\nWe shall start with the spinless case, where at $T=0$ and $t\\gg\\hbar\/E_C$\na simple calculation gives\n\\begin{eqnarray*}\n\\langle H'\\!(t)H'\\!(0)\\rangle &=& \\frac{\\hbar\\Gamma_0(N)}{2\\pi i t}\n\\\\\n &&+ \\frac{\\gamma\\hbar^3(|r_L|^2 + |r_R|^2 - 2|r_L||r_R|\\cos2\\pi N)}\n {4\\pi i E_Ct^3},\n\\end{eqnarray*}\nwith $\\Gamma_0(N)$ defined by Eq.~(\\ref{Gamma_0}). The long-time\nasymptotics of the correlator $\\langle H'(t)H'(0)\\rangle$ is given by the\nfirst term, which is proportional to $t^{-1}$, unless $\\Gamma_0=0$. Thus\nthe dimension of the perturbation $H'$ is $\\frac12<1$, and the\nstrong-coupling fixed point is unstable. At a half-integer $N$ and in a\nperfectly symmetric case, $|r_L|=|r_R|$, the parameter $\\Gamma_0$\nvanishes, and the dimension of the operator $H'$ becomes $\\frac32>1$. In\nthis case the strong-coupling fixed point is stable. In the language of\nthe corresponding 2-channel Kondo problem this means that the\nstrong-coupling fixed point is stable only in the absence of magnetic\nfield and channel anisotropy. The leading irrelevant operator with\ndimension $\\frac32$ gives rise to the weak temperature dependence of the\npeak heights in the symmetric spinless case, see Appendix~\\ref{A1}.\n\nA similar analysis can be performed in the case of electrons with spins,\ndescribed by the Hamiltonian\n(\\ref{H_0_with_spins})--(\\ref{backscattering_with_spins}). The correlator\nof the perturbation operators now has the form\n\\[\n\\langle H'(t)H'(0)\\rangle =\n \\bigl(|r_L|^2+|r_R|^2\\bigr)\n \\sqrt{\\frac{4\\gamma E_C}{\\pi^5}}\\,\n \\left(\\frac{\\hbar}{it}\\right)^{3\/2},\n\\]\nmeaning that $H'$ has dimension $\\frac34<1$ at any $N$ and any barrier\nstrengths. Thus as expected the strong-coupling fixed point in the\n4-channel case is unstable. At a half-integer $N$ this statement has a\nwell-known analog in the theory\\cite{NB} of the multichannel Kondo\nproblem, where the strong-coupling fixed point is unstable even in\nthe absence of the magnetic field.\n\nThe fact that the strong-coupling fixed point is unstable means that as\nthe temperature is lowered, the weak backscattering in the contacts\nbecomes stronger. It is natural to assume that the effective barriers grow\nindefinitely, and the system approaches the weak-coupling fixed point. To\ntest this hypothesis we will consider a system where the two contacts are\nin the weak-tunneling regime and study the stability of the corresponding\nfixed point.\n\nAs we saw in Sec.~\\ref{weak} and \\ref{withspin}, in the weak-tunneling\ncase it is convenient to describe the system in terms of the original\nfermionic variables. Generalizing the Hamiltonian of Sec.~\\ref{asymmetric}\nto the case of both barriers being in the weak-tunneling regime, we\nintroduce the Hamiltonian $H=H_0 + H_C + H_L'+H_R'$, where\n\\begin{eqnarray}\nH_0 &=& \\sum_{k\\alpha} \\epsilon_k a_{k\\alpha}^\\dagger a_{k\\alpha}^{}\n + \\sum_{p\\alpha} \\epsilon_p a_{p\\alpha}^\\dagger a_{p\\alpha}^{}\n\\nonumber\\\\\n & & + \\sum_{q\\alpha} \\epsilon_q a_{q\\alpha}^\\dagger a_{q\\alpha}^{}\n + \\sum_{s\\alpha} \\epsilon_s a_{s\\alpha}^\\dagger a_{s\\alpha}^{} ,\n\\label{H_0_both_weak}\\\\\nH_C &=& E_C(\\hat n_L + \\hat n_R - N)^2,\n\\label{H_C_both_weak}\\\\\n H_L' &=& \\sum_{kp\\alpha}\n (v_L^{} a_{k\\alpha}^\\dagger a_{p\\alpha}^{}F_L^{}\n + v_L^* a_{p\\alpha}^\\dagger a_{k\\alpha}^{}F_L^\\dagger),\n\\label{H_L'_both_weak}\\\\\n H_R' &=& \\sum_{qs\\alpha}\n (v_R^{} a_{s\\alpha}^\\dagger a_{q\\alpha}^{}F_R^{}\n + v_R^* a_{q\\alpha}^\\dagger a_{s\\alpha}^{}F_R^\\dagger).\n\\label{H_R'_both_weak}\n\\end{eqnarray}\nHere $a_{k\\alpha}$ and $a_{s\\alpha}$ describe the electrons in the left\nand right leads, respectively; $a_{p\\alpha}$ and $a_{q\\alpha}$ correspond\nto the electrons in the dot coupled to the left and right contacts.\nSimilarly to the formalism of Sec.~\\ref{asymmetric}, we have introduced\nthe operators $\\hat n_L$ and $\\hat n_R$ of the number of electrons\ntunneled through the left and right contacts and the corresponding raising\nand lowering operators $F_L^{}$, $F_L^\\dagger$, $F_R^{}$, and\n$F_R^\\dagger$ in the tunnel Hamiltonians (\\ref{H_L'_both_weak}) and\n(\\ref{H_R'_both_weak}).\n\nAway from the resonance values of the gate voltage, $N\\neq \\pm\\frac12,\n\\pm\\frac32,\\ldots$, the states with different $n_L$ and $n_R$ have different\nenergies because of the Coulomb term (\\ref{H_C_both_weak}). Since\nthe tunnel Hamiltonian given by (\\ref{H_L'_both_weak}) and\n(\\ref{H_R'_both_weak}) changes $n_L$ or $n_R$, it does not describe the\ntransitions between the low-energy states of the Hamiltonian. Thus when\nconsidering the low-energy properties of the system, one should use the\ntunneling terms (\\ref{H_L'_both_weak}) and (\\ref{H_R'_both_weak}) to\nconstruct the perturbation which is not accompanied by the change of the\ntotal charge of the dot. In the simplest case, this can be done by\nselecting the terms containing equal number of operators $F_i$ and\n$F_i^\\dagger$ in the second-order perturbation\n\\begin{equation}\n H_2' = (H_L' + H_R')\\frac{1}{E_0 - H_0 - H_C}(H_L' + H_R').\n\\label{H_2'}\n\\end{equation}\nHere the denominator accounts for the energy of the virtual state created\nby adding (or removing) an electron to the dot. At very low energies the\nonly important contribution to the denominator is the change of the\nelectrostatic energy of the system.\n\nAt $N$ in the interval $\\bigl(-\\frac12, \\frac12\\bigr)$ a typical\nrepresentative of the variety of terms contained in Eq.~(\\ref{H_2'}) has\nthe form\n\\begin{eqnarray}\nA &=& -v_L^* v_R^{}\\left[\\frac{1}{E_C(1-2N)}+\\frac{1}{E_C(1+2N)}\\right]\n\\nonumber\\\\\n & & \\times\n F_L^\\dagger F_R^{}\n \\sum_{kp\\alpha} a_{p\\alpha}^\\dagger a_{k\\alpha}^{}\n \\sum_{qs\\alpha'} a_{s\\alpha'}^\\dagger a_{q\\alpha'}^{}.\n \\label{A}\n\\end{eqnarray}\nThe terms included in Eq.~(\\ref{A}) correspond to the processes where one\nelectron tunnels into the dot through the left junction and another one\nescapes from the dot through the right junction. Depending on which of the\ntwo tunneling events is executed first, we get different energies of the\nvirtual states, corresponding to the two denominators in Eq.~(\\ref{A}).\n\nTo find out whether or not the weak-coupling fixed point is stable, one\nshould evaluate the scaling dimension of operator $A$ and other similar\ncontributions in Eq.~(\\ref{H_2'}). A straightforward calculation gives\n\\begin{eqnarray}\n\\langle A(t) A^\\dagger(0)\\rangle &=&\n \\frac{\\hbar^2G_L G_R}{(2\\pi e^2)^2}\\!\n \\left[\\frac{1}{E_C(1-2N)}+\\frac{1}{E_C(1+2N)}\\right]^2\n\\nonumber\\\\\n & &\\times\n \\left\\{\\frac{\\pi T}{\\sinh[\\pi T(t-i\\delta)\/\\hbar]}\\right\\}^4.\n \\label{A_correlator}\n\\end{eqnarray}\nAt $T\\to0$ the correlator decays as $1\/t^4$, and therefore\nthe dimension of the operator $A$ is 2. One can also check that the other\ncontributions to Eq.~(\\ref{H_2'}) have the same dimension, whereas the\nhigher-order terms have even higher dimensions. Thus the weak-coupling\nfixed point is stable.\n\nSo far we have shown that away from the resonance values of the gate\nvoltage, $N\\neq \\pm\\frac12, \\pm\\frac32, \\ldots$, the strong-coupling fixed\npoint is unstable whereas the weak-coupling fixed point is stable. We now\nconjecture that there are no fixed points at intermediate coupling. Then\nat low energies the system always evolves towards the stable\n(weak-coupling) fixed point. One should also expect that near this fixed\npoint the temperature dependence of conductance is determined by the\ndimension of the leading irrelevant perturbation, which is a universal\ncharacteristic of the problem. Indeed, one can easily show that the\nleading contribution to the current operator $I=e\\dot n_L$ is given by\n\\begin{equation}\nI = -\\frac{ie}{\\hbar}(A-A^\\dagger).\n \\label{current_in_terms_of_A}\n\\end{equation}\nWe now substitute this current operator into the Kubo formula\n(\\ref{Kubo}), and use (\\ref{A_correlator}) to find the following\nexpression for the conductance:\n\\begin{equation}\nG = \\frac{\\pi\\hbar G_LG_R}{3e^2}\n \\left[\\frac{1}{E_C(1-2N)}+\\frac{1}{E_C(1+2N)}\\right]^2 T^2.\n \\label{cotunneling_conductance}\n\\end{equation}\nThis expression for the conductance in the weak-tunneling limit coincides\nwith the well-known result for the inelastic co-tunneling,\\cite{AN} and\nreproduces Eq.~(\\ref{cotunneling}) near the resonances,\n\\mbox{$\\frac12\\pm N \\ll 1$}.\n\nIt is important to note that in the derivation of\nEq.~(\\ref{cotunneling_conductance}) we neglected all the contributions\nto the current operator with dimensions $d>2$. Such terms would lead to\nthe additional terms in the correlator (\\ref{A_correlator}), which are\nproportional to $\\{\\pi T\/\\sinh[\\pi T(t-i\\delta)\/\\hbar]\\}^m$ with $m>4$.\nClearly, the corresponding corrections to the conductance would behave as\n$\\delta G \\propto T^{m-2}\\ll T^2$. Thus at $T\\to 0$ the temperature\ndependence of the conductance is given by the leading irrelevant\nperturbation. As we have seen, the universal scaling dimension $d=2$ of\nthis perturbation translates into the universal temperature dependence of\nthe conductance $G\\propto T^2$. Therefore this temperature dependence is\nspecific not only for the cases of weak or strong tunneling, but should\nhold for any barrier strengths.\n\nThe above treatment of the stability of the weak-coupling fixed point was\nlimited to the off-resonance values of the gate voltage. We will now\ndemonstrate that the properties of the system are completely different at\nhalf-integer $N$. Consider, for example, the case of $N=\\frac12$. Then the\nstates with the dot charge equal to 0 and $e$ have the same energy. As a\nresult the terms in the tunnel Hamiltonians (\\ref{H_L'_both_weak}) and\n(\\ref{H_R'_both_weak}) responsible for the transitions between these two\nstates no longer lead to the large increase in the energy of the system\nand should not be discounted. It is easy to check that the dimension of\nsuch operators is $d=1$, i.e., the tunnel Hamiltonian is a marginal\noperator. This conclusion is actually obvious, because the same\ncalculation of scaling dimension must be valid in the absence of the\ncharging effects, $E_C=0$, where we do not expect the tunneling amplitudes\nto be renormalized as the temperature is lowered. However, the simple\ncalculation of the scaling dimension is not sufficient to determine the\nstability of the weak-coupling fixed point. To make such a determination,\none has to study the scaling equations for the tunneling amplitudes\nincluding the higher-order terms in $v_L$ and $v_R$. We have performed\nthis analysis in Sec.~\\ref{weak}, where we saw from the mapping to the\nmultichannel Kondo problem that the tunneling amplitude is in fact a\nmarginally relevant perturbation, i.e., it grows logarithmically at low\ntemperatures. Therefore on resonance the weak-coupling fixed point is\nunstable.\n\nThe on-resonance scaling properties of our system are strongly affected by\nthe presence of electron spins and by the symmetry of the two barriers.\nLet us first consider the case of symmetric barriers. In the spinless\ncase we saw that on resonance the strong-coupling fixed point is stable if\nthe barriers are symmetric. Thus one should expect the system to always\nrenormalize towards this fixed point. Then the conductance must be on the\norder of $e^2\/\\hbar$. This is confirmed by the explicit calculation in the\nstrong-tunneling case, Eq.~(\\ref{conductance_spinless}).\n\nIn the presence of spins, both weak- and strong-coupling fixed points are\nunstable on resonance. Thus one should expect that there must be a stable\nfixed point at some intermediate coupling strength. This hypothesis is\nalso confirmed by the analogy with the 4-channel Kondo problem discussed\nin Sec.~\\ref{weak}. Indeed, it is well-known that the low-energy\nproperties of the Kondo problem with the number of channels exceeding 2 is\ngoverned by an intermediate-coupling fixed point, which gives rise to a\nspecific non-Fermi-liquid behavior of the susceptibility and specific\nheat.\\cite{NB,bethe,Affleck,Fabrizio} In Appendix~\\ref{4-channel} we show\nhow one can study the properties of the intermediate-coupling fixed point\nstarting from the vicinity of the strong-coupling point. We discover that\non resonance the thermodynamics of our problem is indeed identical to that\nof the 4-channel Kondo problem.\n\nThe technique presented in Appendix~\\ref{4-channel}, however, does not\nenable us to find the low-temperature conductance $G$. Nevertheless one\ncan estimate the order of magnitude of $G$ at the intermediate-coupling\nfixed point in the following way. From the fact that the strong-coupling\nfixed point is unstable and from the result (\\ref{2nd_order_with_spins})\nof the perturbation theory near this point we know that the conductance\nmust be smaller than $e^2\/2\\pi\\hbar$. On the other hand, at $G\\ll\ne^2\/\\hbar$ the weak-coupling treatment of Sec.~\\ref{weak} is applicable\nand predicts the growth of conductance at $T\\to0$. Therefore the limiting\nvalue of conductance must be of the order, but smaller than\n$e^2\/2\\pi\\hbar$. We further conjecture that the peak conductance is a\nuniversal characteristic of the intermediate-coupling fixed point, see\nAppendix~\\ref{4-channel}, meaning that in the limit $T\\to0$ the peak\nconductance does not depend on the heights of the barriers.\n\nThe scaling analysis in the spirit of Ref.~\\onlinecite{Kane} enables us to\nmake some conclusions about the shape of the peaks in $G(N)$. In the\nspinless case we see from Eq.~(\\ref{conductance_spinless}) that at the\ntails of the peak at $N=\\frac12$ the conductance decays as $G\\propto\n1\/\\bigl(N-\\frac12\\bigr)^4$. This result can be interpreted within the\nscaling picture as follows. A small deviation of $N$ from $\\frac12$ is\nidentified as a magnetic field $h$ in the appropriate multichannel Kondo\nproblem, see Eq.~(\\ref{Kondo}). Therefore such a deviation is a relevant\nperturbation with scaling dimension $d=2\/(2+k)$, where $k$ is the number\nof channels.\\cite{Affleck} In the spinless case we have $k=2$ and\n$d=\\frac12$. This means that as the temperature is lowered, the magnetic\nfield grows as $h(T)\\propto \\bigl(N-\\frac12\\bigr)T^{-1\/2}$. We now\nconjecture that near the fixed point, at $T\\to0$, the conductance is a\nuniversal function of $h$ only, $G=G(h(T))$. As we discussed above, away\nfrom the center of the peak $G\\propto T^2$ at $T\\to0$. Therefore the\nmagnetic field dependence of the conductance must be given by $G\\propto\n1\/h^4\\propto T^2\/\\bigl(N-\\frac12\\bigr)^4$, in agreement with\nEq.~(\\ref{conductance_spinless}). We can now apply the same argument to\nthe case of electrons with spins, where the number of channels $k=4$ and\nthe scaling dimension of the magnetic field is $d=\\frac13$. The\ntemperature dependence of the renormalized magnetic field is consequently\n$h(T) \\propto \\bigl(N-\\frac12\\bigr)T^{-2\/3}$. Therefore near the\nintermediate-coupling fixed point the condition $G\\propto T^2$ at $T\\to 0$\ndemands $G\\propto h^{-3}\\propto T^2\/\\bigl(N-\\frac12\\bigr)^3$. Thus we\nexpect the tails of the conductance peaks to decay as\n$1\/\\bigl(N-\\frac12\\bigr)^3$.\n\n\nSo far we discussed the case of symmetric barriers. The on-resonance\nbehavior of the system is very different if the two barriers are not\nidentical. This can be readily seem from the analogy to the Kondo problem.\n{}From the theory of the multichannel Kondo model\\cite{NB,Melnikov} it is\nknown that the asymmetry between the channels is a relevant perturbation.\nIn terms of the Coulomb blockade problem this means that on resonance a\nsmall asymmetry will grow and eventually the smaller transmission\ncoefficient, say ${\\cal T}_L$, will scale to zero, whereas the larger one,\n${\\cal T}_R$, must approach the fixed point of the problem with half the\nnumber of channels of the original one. Thus for both 2- and 4-channel\ncases ${\\cal T}_R \\to 1$. For the most realistic case of electrons with\nspins, near this asymmetric fixed point the system is identical to the one\nconsidered in Sec.~\\ref{asymmetric}. Thus we conclude that for asymmetric\nsystems (i) the on-resonance conductance vanishes linearly in $T$, and (ii)\nthe shape of the resonances in $G(N)$ is universal and given by\nEq.~(\\ref{conductance_asymmetric_final}).\n\nThe on-resonance scaling properties of our problem are summarized in\nFig.~\\ref{fig:4}.\n\n\\vfill\n\\begin{figure}\n\\narrowtext\n\\epsfxsize=2.6in\n\\hspace{.1em}\n\\epsffile{4chan4.ps}\n\\vspace{1.5\\baselineskip}\n\\caption{Scaling diagram of the two-barrier system in the cases of (a)\n spinless electrons and (b) electrons with spins. The gate voltage\n corresponds to a resonance value, $N=\\pm\\frac12,\\pm\\frac32,\\ldots$. In\n the symmetric case, ${\\cal T}_L={\\cal T}_R$, the system scales to either\n strong- or intermediate-coupling fixed point, depending on the number of\n channels. In an asymmetric case the system approaches the limit where\n the weaker barrier disappears, while the stronger one becomes very\n high.}\n\\label{fig:4}\n\\end{figure}\n\n\n\\section{Conclusion}\n\nIn this paper we have studied the transport through a quantum dot\nconnected to two leads by quantum point contacts, Fig.~\\ref{fig:1}. Unlike\nin conventional theories of Coulomb blockade,\\cite{Averin,Grabert} we\nconcentrate on the regime of strong tunneling, where one or both contacts\nare close to perfect transmission. We find the linear conductance in a\nwide interval of temperatures, limited by the level spacing in the dot\nfrom below and by the charging energy $E_C$ from above.\n\nWe have shown that as long as there is any backscattering in the contacts,\ni.e., both transmission coefficients ${\\cal T}_{L,R}<1$, the Coulomb\nblockade peaks in the conductance as a function of the gate voltage should\nbe observed. Between the Coulomb blockade peaks the temperature dependence\nof the conductance at $T\\to0$ is determined by the weak-coupling fixed\npoint. As a result we get the $T^2$ dependence of the conductance, which\ncoincides with the known result\\cite{AN} for the inelastic co-tunneling\nmechanism.\n\nOur results show that the behavior of the conductance near the peaks\ndepends qualitatively on the presence of electron spins. In the spinless\ncase the strong-tunneling fixed point is stable if the barriers are\nsymmetric, and thus the technique developed in this paper allows a\ncomplete solution of the problem. The resulting peaks in conductance have\nthe height $e^2\/4\\pi\\hbar$, with the shape given by\nEq.~(\\ref{conductance_spinless}). In the presence of the spins, the\nstrong-tunneling fixed point is no longer stable, and our technique does\nnot allow a complete solution of the problem. Using a useful analogy with\nthe multichannel Kondo problem, we have concluded that the peak\nconductance must still be of the order of $e^2\/\\hbar$. In addition, we\nstudied the peaks in conductance in the case of asymmetric\nbarriers. The peak conductance in this case is linear in temperature at\n$T\\to0$, and its shape is given by\nEq.~(\\ref{conductance_asymmetric_final}).\n\nThe strong-tunneling regime of Coulomb blockade was explored to some\nextent in several recent experiments.\\cite{van,Pasquier,Waugh,Molenkamp}\nOur results are in qualitative agreement with the experiment.\\cite{van} We\nbelieve that our predictions can be easily tested in similar experiments.\n\n\n\n\n\n\\acknowledgements\n\nWe are grateful to L.~I.~Glazman and P.~A.~Lee for helpful discussions.\nThe work was sponsored by Joint Services Electronics Program Contract\nDAAH04-95-1-0038.\n\n\\end{multicols}\n\\widetext\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction} \n\nLet $G \\subset GL_n(\\mathbb{C})$ be a finite complex reflection group and denote by ${\\mathcal A}(G)$ the union of all the reflecting hyperplerplanes of $G$, i.e. the hyperplanes in $\\mathbb{C}^n$ fixed by some $g \\in G$, $g \\ne Id$. For a general reference on complex reflection groups, see \\cite{LT} and \\cite{OT}.\nSince ${\\mathcal A}(G)$ is a central hyperplane arrangement, it has a defining equation $f=0$ in $\\mathbb{C}^n$, where $f$ \nis a homogeneous polynomial of some degree $d$. One can associate to this setting the Milnor fiber $F(G)$ of the arrangement ${\\mathcal A}(G)$. This is a smooth hypersurface in $\\mathbb{C}^n$, defined by $f=1$, and it is endowed with a monodromy morphism $h: F(G) \\to F(G)$, given by $h(x_1,...,x_n)=\\exp(2\\pi i\/d)\\cdot (x_1,...,x_n)$, see \\cite{BDY, BaSe, BDS, CS, DHA, DL1, DL2, L, MP, Se1, Se2, S2} for related results and to get a feeling of the problems in this very active area.\nThe study of the induced monodromy operator\n$$h^1(G): H^1(F(G),\\mathbb{C}) \\to H^1(F(G),\\mathbb{C})$$\nis the object of the papers \\cite{MPP} and \\cite{Dreflection}, while in the special case of real reflection groups $G$, there are some additional results on the higher degree monodromy operators \n$$h^j(G): H^j(F(G),\\mathbb{C}) \\to H^j(F(G),\\mathbb{C})$$\nwhere $j>1$, see \\cite{Se1, Se2, DL2}. In general, not only these monodromy operators are not known, but even the Betti numbers $b_j(F(G))$ are known only in a limited number of cases.\n\nThe complex reflection groups have been classified by Shephard and Todd, see \\cite{ST}, who showed that there is an infinite series $G(m,p,n)$ of such groups, plus 34 exceptional cases.\nThe exceptional complex reflection groups in this classification are usually denoted by $G_j$, with $4 \\leq j \\leq 37$. \nConsider the following polynomial of degree 60 in $S=\\mathbb{C}[x,y,z,t]$ \n\\begin{equation}\n\\label{e6}\nf=xyzt(x^4-y^4)(x^4-z^4)(x^4-t^4)(y^4-z^4)(y^4-t^4)(z^4-t^4)\n\\end{equation}\n$$((x-y)^2-(z+t)^2) ((x-y)^2-(z-t)^2) ((x+y)^2-(z+t)^2)((x+y)^2-(z-t)^2)$$\n$$((x-y)^2+(z+t)^2)((x-y)^2+(z-t)^2)((x+y)^2+(z+t)^2)((x+y)^2+(z-t)^2)$$\n$$((x-z)^2+(y+t)^2) ((x-z)^2+(y-t)^2)((x+z)^2+(y+t)^2)((x+z)^2+(y-t)^2)$$\n$$((x-t)^2+(y+z)^2)((x-t)^2+(y-z)^2)((x+t)^2+(y+z)^2)((x+t)^2+(y-z)^2).$$\nThen the reflection hyperplane arrangement ${\\mathcal A}(G_{31})$ consists of 60 hyperplanes in $\\mathbb{C}^4$ and is given by $f=0$, see \\cite{HR}.\nIn this note we prove the following result.\n\n\\begin{thm}\n\\label{thmmain}\nThe monodromy operator\n$$h^1:H^1(F(G_{31}),\\mathbb{C}) \\to H^1(F(G_{31}),\\mathbb{C})$$\nfor the exceptional complex reflection group $G_{31}$ is the identity. In particular, the first Betti number $b_1(F(G_{31}))$ is $59$.\n\\end{thm}\nThe description of the monodromy operator\n$h^1$ for all the other complex reflection groups is given in \\cite{Dreflection}, using a method which cannot be applied to the exceptional group $G=G_{31}$ as explained in \n\\cite[Remark 6.2]{Dreflection}. The proof for the case $G=G_{31}$ involves a completely different approach, and this explains why it is written down here as a separate note. Indeed, this proof is close in spirit to our paper \\cite{DStCompM}, with computer aided computations playing a key role at several stages. Moreover, the rank of the group $G_{31}$ being $n=4$, we have to deal with higher dimensional singularities and hence with a more complicated spectral sequence than in \\cite{DStCompM}, where the case of plane curves corresponding to $n=3$ is discussed, see also Remark \\ref{rkdim} (i).\n\nIn the second section we recall a spectral sequence approach for the computation of the monodromy of the Milnor fibers of homogeneous polynomials with arbitrary singularities introduced in \\cite{Dcomp}, and developed in \\cite[Chapter 6]{D1}, \\cite{DStCompM}, with several key additions in the joint work of Morihiko Saito and the first author, see \\cite{DS1}.\nFor simplicity, we describe the results only in the case $n=4$, the only case needed in the sequel. However, the approach presented here is very general, and can be applied at least to all free hyperplane arrangements to get valuable information on their first Milnor cohomology.\nOn the other hand, the success of our method is based on the special properties enjoyed by the mixed Hodge structure on the first cohomology group $H^1(F(G),\\mathbb{Q})$, see the proof of Theorem \\ref{thmEV}, and hence this simple approach does not give complete results on\nthe higher cohomology groups $H^j(F(G),\\mathbb{Q})$, where $j>1$.\n\n\nIn the third section we recall basic facts on the exceptional complex reflection group $G_{31}$, in particular the construction of the basic invariants going back to Maschke \\cite{Mas} and the construction of a basis of Jacobian syzygies as described by Orlik and Terao in \\cite[Appendix B, pp. 280-281 and p. 285]{OT}.\n\nIn the fourth section we describe the algorithm used to determine the monodromy operator\n$h^1:H^1(F(G_{31}),\\mathbb{C}) \\to H^1(F(G_{31}),\\mathbb{C})$. For this we need a careful study of the second cohomology group of the Koszul complex of the partial derivatives $f_x$, $f_y$, $f_z$ and $f_t$ of the polynomial $f$ of degree 60 given in \\eqref{e6}. At several points, this study is done by\nusing the software SINGULAR \\cite{Sing}. To get an idea why computations via SINGULAR are necessary, let us remark that each of the sequences of three homogeneous polynomials in $S$ which are shown to be regular sequences in Proposition \\ref{prop1} consists of polynomials of degrees\n$28, 12$ and $16$ respectively, and having $136,24$ and respectively $45$ monomials occurring with non zero (and, in fact, quite complicated rational) coefficients. And the final step in the proof is Lemma \\ref{lemfinal} where a homogeneous system of 19600 linear equations in 1424 indeterminates is shown to admit only the trivial solution.\nThe corresponding codes are available at \\url{http:\/\/math.unice.fr\/~dimca\/singular.html}.\n\n\\medskip\n\nWe would like to thank the referee for his useful suggestions.\n\n\\section{A spectral sequence for the Milnor monodromy} \nLet $S=\\mathbb{C}[x,y,z,t]$ and let $f \\in S$ be a homogeneous polynomial of degree $d$ without multiple factors. Consider the global Milnor fiber $F$ of the polynomial $f$ defined by $f(x,y,z,t)=1$ in $\\mathbb{C}^4$, with monodromy action $h:F \\to F$, \n$$h(x,y,z,t)=\\exp(2\\pi i\/d)\\cdot (x,y,z,t).$$\n\nLet $\\Omega^j$ denote the graded $S$-module of (polynomial) differential $j$-forms on $\\mathbb{C}^4$, for $0 \\leq j \\leq 4$. The complex $K^*_f=(\\Omega^*, d f \\wedge)$ is nothing else but the Koszul complex in $S$ of the partial derivatives $f_x$, $f_y$, $f_z$ and $f_t$ of the polynomial $f$. The general theory says that there is a spectral sequence $E_*(f)$, whose first term $E_1(f)$ is computable in terms of the homogeneous components of the cohomology of the Koszul complex $K^*_f$ and whose limit\n$E_{\\infty}(f)$ gives us the action of monodromy operator on the graded pieces of $H^*(F,\\mathbb{C})$ with respect to a certain pole order filtration $P$, see for details \\cite{Dcomp}, \\cite[Chapter 6]{D1}, \\cite{DS1}.\n\n More precisely, for any integer $k \\in [1,d]$, there is a spectral sequence \n\\begin{equation} \n\\label{newspsq}\nE_1^{s,t}(f)_k=H^{s+t+1}(K^*_f)_{td+k},\n\\end{equation} \nconverging to\n\\begin{equation} \n\\label{limit}\nE_{\\infty}^{s,t}(f)_k=Gr_P^sH^{s+t}(F,\\mathbb{C})_{\\lambda },\n\\end{equation} \nwhere $\\lambda=\\exp (-2 \\pi ik\/d)$ and $H^{s+t}(F,\\mathbb{C})_{\\lambda }$ denotes the associated eigenspace of the monodromy operator.\nMoreover, the differential $d_1$ in this spectral sequence is induced by the exterior differentiation of forms, i.e. $d_1 :[\\omega] \\mapsto [d (\\omega)]$. We have the following.\n\n\\begin{thm}\n\\label{thmEV}\nLet $D:f=0$ be a reduced degree $d$ surface in $\\mathbb{P}^3$, and let $$\\lambda=\\exp (-2\\pi i k\/d)\\ne 1,$$ with $k \\in (0,d)$ an integer. \nThen $H^1(F,\\mathbb{C})_{\\lambda}=0$ if and only if\n$$E_2^{1,0}(f)_k =E_2^{1,0}(f)_{k'} = 0,$$\nwhere $k'=d-k$. \n\\end{thm}\n\\proof\nNote that on the cohomology group $H^1(F,\\mathbb{C})_{\\ne 1}$, which is by definition the direct sum of eigenspaces of $h^1$ corresponding to eigenvalues $\\ne 1$, one has $P^2=0$, $F^1 \\subset P^1$\nand $F^0 = P^0=H^1(F,\\mathbb{C})_{\\ne 1}$, where $F^1$ denotes the Hodge filtration, see \\cite{DS1}. Let now $\\lambda$ be an eigenvalue of the monodromy operator.\nThe fact that $H^1(F,\\mathbb{C})_{\\ne 1}$ is a pure Hodge structure of weight 1, see \\cite{BDS}, \\cite{DP}, implies that either\n$H^{1,0}(F,\\mathbb{C})_{\\lambda } \\ne 0$ or $H^{0,1}(F,\\mathbb{C})_{\\lambda } \\ne 0$. In the first case we get\n$Gr_P^{1}H^{1}(F,\\mathbb{C})_{\\lambda }\\ne 0$, while in the second case we get $H^{1,0}(F,\\mathbb{C})_{\\overline \\lambda } \\ne 0$ and hence $Gr_P^{1}H^{1}(F,\\mathbb{C})_{\\overline \\lambda }\\ne 0$.\nThe result follows now from the fact that $H^1(K^*_f)=0$ since $f$ is a reduced polynomial.\nIndeed, it is known that the smallest $j$ with $H^j(K^*_f) \\ne 0$ is precisely the codimension of the singular locus $D_{sing}$ of $D$ in $\\mathbb{P}^3$, see \\cite[Theorems A.2.38 and A.2.48]{Eis}. This implies that $E_1^{s,t}(f)_k=0$ for $s+t=0$, and hence \n$$E_2^{1,0}(f)_k =E_{\\infty}^{1,0}(f)_k$$\nfor any $k \\in (0,d)$.\n\\endproof\nTo check the condition $E_2^{1,0}(f)_k =0$ in practice, we proceed as follows. Consider a 2-form\n \\begin{equation}\n\\label{e9}\n\\omega= \\sum_{1\\leq i2} $ is zero, see Lemma \\ref{p3 zero average}. \nActually we only need that the average of \nthe function in \\eqref{unico pezzo} is zero. If \n $ f_5 = 0 $ \n (see \\eqref{order5}) then $({\\mathtt S}1)$ is not required. \nThis completes the task of conjugating ${\\cal L}_\\omega $ to $ {\\cal L}_6 $ in \\eqref{L6 qualitativo}.\n\nFinally, in section \\ref{subsec:mL0 mL5} we apply the abstract reducibility Theorem 4.2 in \\cite{BBM}, \nbased on a quadratic KAM scheme, which \ncompletely diagonalizes the linearized operator, obtaining \\eqref{linearosc}. \nThe required smallness condition \\eqref{R6resto} for $ R_6 $ holds. Indeed\nthe biggest term in $ R_6 $ comes from the conjugation of \n$ \\varepsilon \\partial_x v_\\varepsilon (\\theta_0 (\\varphi), y_\\d (\\varphi)) $ in \\eqref{a1p1p2}.\nThe linear BNF procedure of section \\ref{BNF:step1} had eliminated \nits main contribution $ \\varepsilon \\partial_x v_\\varepsilon ( \\varphi,0) $. \nIt remains $ \\varepsilon \\partial_x \\big( v_\\varepsilon (\\theta_0 (\\varphi), y_\\d (\\varphi) ) - v_\\varepsilon ( \\varphi,0) \\big) $ \nwhich has size $ O( \\varepsilon^{7-2b} \\gamma^{-1} ) $ due to the estimate \\eqref{ansatz 0} of the approximate solution.\nThis term enters in the variable coefficients of $ d_1 (\\varphi,x) \\partial_x $\nand $ d_0 (\\varphi, x) \\partial_x^0 $. The first one had been reduced to the constant \noperator $ m_1 \\partial_x $ by the descent method of section \\ref{step5}.\nThe latter term is an operator of order $O(\\partial_x^0 )$ which satisfies \\eqref{R6resto}.\nThus $ {\\cal L}_6 $ may be diagonalized \nby the iterative scheme of Theorem 4.2 in \\cite{BBM} \nwhich requires the smallness condition $ O( \\varepsilon^{7-2b} \\gamma^{-2}) \\ll 1 $. This is the content of section \\ref{subsec:mL0 mL5}. \n\\\\[1mm]\n{\\it The Nash-Moser iteration.}\nIn section \\ref{sec:NM} we perform the nonlinear Nash-Moser iteration which \nfinally proves Theorem \\ref{main theorem} and, therefore, Theorem \\ref{thm:KdV}.\nThe optimal smallness condition required for the convergence of the scheme is \n$ \\varepsilon \\| {\\cal F}(\\varphi, 0, 0 ) \\|_{s_0+ \\mu} \\gamma^{-2} \\ll 1 $, see \\eqref{nash moser smallness condition}.\nIt is verified because $ \\| X_P(\\varphi, 0 , 0 ) \\|_s \\leq_s \\varepsilon^{6 - 2b} $ \n(see \\eqref{stima XP}), which, in turn, is a consequence of having eliminated the terms\n$ O(v^5), O( v^4 z)$ from the original Hamiltonian \\eqref{H iniziale KdV}, \nsee \\eqref{widetilde cal H}. This requires the condition ($ {\\mathtt S}2$). \n\n\\smallskip\n\n\\noindent\n{\\it Acknowledgements}. We thank M. Procesi, P. Bolle and T. Kappeler for many useful discussions. \nThis research was supported by the European Research Council under FP7, \nand partially by the grants STAR 2013 and \nPRIN 2012 ``Variational and perturbative aspects of nonlinear differential problems\".\n\n\\section{Preliminaries}\n\n\\subsection{Hamiltonian formalism of KdV}\\label{sec:Ham For}\n\nThe Hamiltonian vector field $X_H$ generated by a Hamiltonian \n$ H : H^1_0(\\mathbb T_x) \\to \\mathbb R$ is $ X_H (u) := \\partial_x \\nabla H (u) $, because\n$$\nd H (u) [ h] \n= ( \\nabla H(u), h )_{L^2(\\mathbb T_x)} = \\Omega ( X_H (u), h ) \\, , \\quad \\forall u, h \\in H^1_0(\\mathbb T_x) \\, , \n$$\nwhere $\\Omega$ is the non-degenerate symplectic form\n\\begin{equation}\\label{2form KdV}\n\\Omega (u, v) := \\int_{\\mathbb T} (\\partial_x^{-1 } u) \\, v \\, dx \\, , \\quad \n\\forall u, v \\in H^1_0 (\\mathbb T_x) \\, , \n\\end{equation}\nand $ \\partial_x^{-1} u $ is the periodic primitive of $ u $ with zero average. \nNote that \n\\begin{equation} \\label{def pi 0} \n\\partial_x \\partial_x^{-1} = \\partial_x^{-1} \\partial_x = \\pi_0 \\, , \\quad \n\\pi_0 (u) := u - \\frac{1}{2\\pi} \\int_{\\mathbb T} u(x) \\, dx \\,. \n\\end{equation}\nA map \nis symplectic if it preserves \nthe 2-form $ \\Omega $.\n\nWe also remind that the Poisson bracket between two functions $ F $, $ G : H^1_0(\\mathbb T_x) \\to \\mathbb R$ is\n\\begin{equation}\\label{Poisson bracket}\n\\{ F (u), G(u) \\} := \\Omega (X_F, X_G ) = \\int_{\\mathbb T} \\nabla F(u) \\partial_x \\nabla G (u) dx \\, . \n\\end{equation}\nThe linearized KdV equation at $ u $ is \n$$\nh_t = \\partial_x \\, (\\partial_u \\nabla H)(u) [h ] = X_{K}(h) \\, , \n$$\nwhere $ X_K $ is the KdV Hamiltonian vector field with quadratic Hamiltonian \n$ K = \\frac12 ((\\partial_u \\nabla H)(u)[h], h)_{L^2(\\mathbb T_x)} $ \n$= \\frac12 (\\partial_{uu} H)(u)[h, h] $. By the Schwartz theorem, the Hessian operator $ A := (\\partial_u \\nabla H)(u) $ is symmetric, namely\n$ A^T = A $, with respect to the $ L^2$-scalar product. \n\n\\smallskip\n\n\\noindent\n{\\bf Dynamical systems formulation.} \nIt is convenient to regard the KdV equation also in the Fourier representation \n\\begin{equation}\\label{Fourier}\nu(x) = {\\mathop \\sum}_{j \\in \\mathbb Z \\setminus \\{0\\} } u_j e^{{\\mathrm i} j x} \\, , \\qquad \nu(x) \\longleftrightarrow u := (u_j)_{j \\in \\mathbb Z \\setminus \\{0\\} } \\, , \\quad u_{-j} = \\overline{u}_j \\, ,\n\\end{equation}\nwhere the Fourier indices $ j \\in \\mathbb Z \\setminus \\{ 0 \\}$ by the definition \\eqref{def phase space} of the phase space\nand $u_{-j} = \\overline{u}_j$ because $u(x)$ is real-valued.\nThe symplectic structure writes\n\\begin{equation}\\label{2form0}\n\\Omega = \\frac12 \\sum_{j \\neq 0} \\frac{1}{{\\mathrm i} j} du_j \\wedge d u_{-j} = \\sum_{j \\geq 1} \n\\frac{1}{{\\mathrm i} j} du_j \\wedge d u_{-j}\n\\, , \\qquad \n\\Omega ( u, v ) = \\sum_{j \\neq 0} \\frac{1}{{\\mathrm i} j} u_j v_{-j} = \\sum_{j \\neq 0} \\frac{1}{{\\mathrm i} j} u_j { \\overline v}_{j} \\, , \n\\end{equation}\nthe Hamiltonian vector field $X_H$ and the Poisson bracket $\\{ F, G \\}$ are\n\\begin{equation}\\label{PoissonBr}\n[X_H (u)]_j = {\\mathrm i} j (\\partial_{u_{-j}} H) (u) \\, , \\ \\forall j \\neq 0 \\, ,\n\\quad \n\\{ F (u), G(u) \\} = - {\\mathop \\sum}_{j \\neq 0} {\\mathrm i} j (\\partial_{u_{-j}} F) (u) (\\partial_{u_j} G) (u) \\, . \\end{equation}\n\\noindent\n\\noindent\n{\\bf Conservation of momentum.} \nA Hamiltonian \n\\begin{equation} \\label{cons mom}\nH(u) = \\sum_{j_1, \\ldots, j_n \\in \\mathbb Z \\setminus \\{ 0 \\} } H_{j_1, \\ldots, j_n} u_{j_1} \\ldots u_{j_n}, \n\\quad \nu(x) = \\sum_{j \\in \\mathbb Z \\setminus \\{ 0 \\} } u_j e^{{\\mathrm i} jx},\n\\end{equation}\nhomogeneous of degree $n$, \\emph{preserves the momentum} if the coefficients \n$H_{j_1, \\ldots, j_n}$ are zero for $j_1 + \\ldots + j_n \\neq 0$, \nso that the sum in \\eqref{cons mom} is restricted to integers \n such that $j_1 + \\ldots + j_n = 0$. \nEquivalently, $H$ preserves the momentum if $\\{ H, M \\} = 0$, where $M$ is the momentum \n$ M(u) := \\int_{\\mathbb T} u^2 dx = $ $ \\sum_{j \\in \\mathbb Z \\setminus \\{0\\}} u_j u_{-j} $. \nThe homogeneous components of degree $ \\leq 5 $ \nof the KdV Hamiltonian $ H $ in \\eqref{Ham in intro} \npreserve the momentum because, by \\eqref{order5}, \nthe homogeneous component $f_5$ of degree 5 does not depend on the space variable $x$. \n\n\\smallskip\n\n\\noindent\n{\\bf Tangential and normal variables.} \nLet $\\bar\\jmath_1, \\ldots, \\bar\\jmath_\\nu \\geq 1$ be $\\nu$ distinct integers, \nand $S^+ := \\{ \\bar\\jmath_1, \\ldots, \\bar\\jmath_\\nu \\}$. \nLet $S$ be the symmetric set in \\eqref{tang sites}, \nand $S^c := \\{ j \\in \\mathbb Z \\setminus \\{ 0 \\} : j \\notin S \\}$ its complementary set in $\\mathbb Z \\setminus \\{ 0 \\}$. \nWe decompose the phase space as \n\\begin{equation}\\label{splitting S-S-bot}\nH^1_0 (\\mathbb T_x) := H_S \\oplus H_S^\\bot \\,, \n\\quad \nH_S := \\mathrm{span}\\{ e^{{\\mathrm i} jx} : \\, j \\in S \\} , \\quad \nH_S^\\bot := \\big\\{ u = \\sum_{j \\in S^c} u_j e^{{\\mathrm i} jx} \\in H^1_0(\\mathbb T_x) \\big\\} ,\n\\end{equation}\nand we denote by $\\Pi_S $, \n$\\Pi_S^\\bot $ the corresponding orthogonal projectors. \nAccordingly we decompose\n\\begin{equation} \\label{u = v + z}\nu = v + z, \\qquad \nv = \\Pi_S u := {\\mathop \\sum}_{j \\in S} u_j \\, e^{{\\mathrm i} jx}, \\quad \nz = \\Pi_S^\\bot u := {\\mathop \\sum}_{j \\in S^c} u_j \\, e^{{\\mathrm i} jx} \\, , \n\\end{equation}\nwhere $ v $ is called the {\\it tangential} variable and $ z $ the {\\it normal} one.\nWe shall sometimes identify $ v \\equiv (v_j)_{j \\in S } $ and $ z \\equiv (z_j)_{j \\in S^c } $. \nThe subspaces $ H_S $ and $ H_S^\\bot $ are {\\it symplectic}. \nThe dynamics of these two components is quite different. \nOn $ H_S $ we shall introduce the action-angle variables, see \\eqref{coordinate azione angolo}.\nThe linear frequencies of oscillations on the tangential sites are \n\\begin{equation}\\label{bar omega}\n\\bar\\omega := (\\bar\\jmath_1^3, \\ldots, \\bar\\jmath_\\nu^3) \\in \\mathbb N^\\nu.\n\\end{equation}\n\n\\subsection{Functional setting}\n\n\n{\\bf Norms.} Along the paper we shall use the notation \n\\begin{equation}\\label{Sobolev coppia}\n\\| u \\|_s := \\| u \\|_{H^s( \\mathbb T^{\\nu + 1})} := \\| u \\|_{H^s_{\\varphi,x} }\n\\end{equation}\nto denote the Sobolev norm of functions $ u = u(\\varphi,x) $ in the Sobolev space $ H^{s} (\\mathbb T^{\\nu + 1} ) $. \nWe shall denote by $ \\| \\ \\|_{H^s_x} $ the Sobolev norm in the phase space of functions $ u := u(x) \\in H^{s} (\\mathbb T ) $.\nMoreover $ \\| \\ \\|_{H^s_\\varphi} $ will denote the Sobolev norm of scalar functions, like \nthe Fourier components $ u_j (\\varphi) $. \n\nWe fix $ s_0 := (\\nu+2) \\slash 2 $ so that $ H^{s_0} (\\mathbb T^{\\nu + 1} ) \\hookrightarrow L^{\\infty} (\\mathbb T^{\\nu + 1} ) $ and the spaces \n$ H^s (\\mathbb T^{\\nu + 1} ) $, $ s > s_0 $, are an algebra. \nAt the end of this section \nwe report interpolation properties of the Sobolev norm that will be currently used along the paper.\nWe shall also denote \n\\begin{align} \\label{HSbot nu}\nH^s_{S^\\bot} (\\mathbb T^{\\nu+1}) \n& := \\big\\{ u \\in H^s(\\mathbb T^{\\nu + 1} ) \\, : \\, u (\\varphi, \\cdot ) \\in H_S^\\bot \\ \n\\forall \\varphi \\in \\mathbb T^\\nu \\big\\} \\,,\n\\\\\n\\label{HS nu}\nH^s_{S} (\\mathbb T^{\\nu+1}) \n& := \\big\\{ u \\in H^s(\\mathbb T^{\\nu + 1} ) \\, : \\, u (\\varphi, \\cdot ) \\in H_{S} \\ \n\\forall \\varphi \\in \\mathbb T^\\nu \\big\\} \\,. \n\\end{align}\nFor a function $u : \\Omega_o \\to E$, $\\omega \\mapsto u(\\omega)$, where $(E, \\| \\ \\|_E)$ is a Banach space and \n$ \\Omega_o $ is a subset of $\\mathbb R^\\nu $, we define the sup-norm and the Lipschitz semi-norm\n\\begin{equation} \\label{def norma sup lip}\n\\| u \\|^{\\sup}_E \n:= \\| u \\|^{\\sup}_{E,\\Omega_o} \n:= \\sup_{ \\omega \\in \\Omega_o } \\| u(\\omega ) \\|_E \\, , \n\\quad\n\\| u \\|^{\\mathrm{lip}}_E \n:= \\| u \\|^{\\mathrm{lip}}_{E,\\Omega_o} \n:= \\sup_{\\omega_1 \\neq \\omega_2 } \n\\frac{ \\| u(\\omega_1) - u(\\omega_2) \\|_E }{ | \\omega_1 - \\omega_2 | }\\,,\n\\end{equation}\nand, for $ \\gamma > 0 $, the Lipschitz norm\n\\begin{equation} \\label{def norma Lipg}\n\\| u \\|^{{\\mathrm{Lip}(\\g)}}_E \n:= \\| u \\|^{{\\mathrm{Lip}(\\g)}}_{E,\\Omega_o}\n:= \\| u \\|^{\\sup}_E + \\gamma \\| u \\|^{\\mathrm{lip}}_E \\, . \n\\end{equation}\nIf $ E = H^s $ we simply denote $ \\| u \\|^{{\\mathrm{Lip}(\\g)}}_{H^s} := \\| u \\|^{{\\mathrm{Lip}(\\g)}}_s $. We shall use the notation \n$$\na \\leq_s b \\quad \\ \\Longleftrightarrow \\quad a \\leq C(s) b \\quad\n\\text{for some constant } C(s) > 0 \\, . \n$$ \n\\noindent\n{\\bf Matrices with off-diagonal decay.}\nA linear operator can be identified, as usual, with its matrix representation. \nWe recall the definition \nof the $ s $-decay norm (introduced in \\cite{BB13JEMS}) of an infinite dimensional matrix. \nThis norm is used in \\cite{BBM} for the KAM reducibility scheme of the linearized operators. \n\\begin{definition}\\label{def:norms}\nThe $s$-decay norm of an infinite dimensional matrix $ A := (A_{i_1}^{i_2} )_{i_1, i_2 \\in \\mathbb Z^b } $, $b \\geq 1$, is \n\\begin{equation} \\label{matrix decay norm}\n\\left| A \\right|_{s}^2 := \n\\sum_{i \\in \\mathbb Z^b} \\left\\langle i \\right\\rangle^{2s} \n\\Big( \\sup_{ \\begin{subarray}{c} i_{1} - i_{2} = i \n\\end{subarray}}\n| A^{i_2}_{i_1}| \\Big)^{2} \\, .\n\\end{equation}\nFor parameter dependent matrices $ A := A(\\omega) $, $\\omega \\in \\Omega_o \\subseteq \\mathbb R^\\nu $, \nthe definitions \\eqref{def norma sup lip} and \\eqref{def norma Lipg} become \n\\begin{equation} \\label{matrix decay norm Lip}\n| A |^{\\sup}_s := \\sup_{ \\omega \\in \\Omega_o } | A(\\omega ) |_s \\, , \n\\quad\n| A |^{\\mathrm{lip}}_s := \\sup_{\\omega_1 \\neq \\omega_2} \n\\frac{ | A(\\omega_1) - A(\\omega_2) |_s }{ | \\omega_1 - \\omega_2 | }\\,,\n\\quad\n| A |^{{\\mathrm{Lip}(\\g)}}_s := | A |^{\\sup}_s + \\gamma | A |^{\\mathrm{lip}}_s \\,. \n\\end{equation}\n\\end{definition}\n\nSuch a norm is modeled on the behavior of matrices representing the multiplication\noperator by a function.\nActually, given a function $ p \\in H^s(\\mathbb T^b) $, the multiplication operator $ h \\mapsto p h $ is represented by the T\\\"oplitz matrix \n$ T_i^{i'} = p_{i - i'} $ and $ |T|_s = \\| p \\|_s $. \nIf $p = p(\\omega )$ is a Lipschitz family of functions, then \n\\begin{equation}\\label{multiplication Lip}\n|T|_s^{\\mathrm{Lip}(\\g)} = \\| p \\|_s^{\\mathrm{Lip}(\\g)}\\,.\n\\end{equation}\nThe $s$-norm satisfies classical algebra and interpolation inequalities, see \\cite{BBM}. \n\n \n\\begin{lemma} \\label{prodest}\nLet $A = A(\\omega)$ and $B = B(\\omega)$ be matrices depending in a Lipschitz way on the parameter $\\omega \\in \n\\Omega_o \\subset \\mathbb R^\\nu $. \nThen for all $s \\geq s_0 > b\/2 $ there are $ C(s) \\geq C(s_0) \\geq 1 $ such that\n\\begin{align} \\label{algebra Lip}\n|A B |_s^{{\\mathrm{Lip}(\\g)}} \n& \\leq C(s) |A|_s^{{\\mathrm{Lip}(\\g)}} |B|_s^{{\\mathrm{Lip}(\\g)}} \\, , \n\\\\\n\\label{interpm Lip}\n|A B|_{s}^{{\\mathrm{Lip}(\\g)}} \n& \\leq C(s) |A|_{s}^{{\\mathrm{Lip}(\\g)}} |B|_{s_0}^{{\\mathrm{Lip}(\\g)}} \n+ C(s_0) |A|_{s_0}^{{\\mathrm{Lip}(\\g)}} |B|_{s}^{{\\mathrm{Lip}(\\g)}} .\n\\end{align}\n\\end{lemma}\nThe $ s $-decay norm controls the Sobolev norm, namely \n\\begin{equation}\\label{interpolazione norme miste}\n\\| A h \\|_s^{\\mathrm{Lip}(\\g)} \n\\leq C(s) \\big(|A|_{s_0}^{\\mathrm{Lip}(\\g)} \\| h \\|_s^{\\mathrm{Lip}(\\g)} + |A|_{s}^{\\mathrm{Lip}(\\g)} \\| h \\|_{s_0}^{\\mathrm{Lip}(\\g)} \\big).\n\\end{equation}\nLet now $ b := \\nu + 1 $. \nAn important sub-algebra is formed by the {\\it T\\\"oplitz in time matrices} defined by\n\\begin{equation}\\label{Topliz matrix}\n A^{(l_2, j_2)}_{(l_1, j_1)} := A^{j_2}_{j_1}(l_1 - l_2 )\\, ,\n\\end{equation}\nwhose decay norm \\eqref{matrix decay norm} is\n\\begin{equation}\\label{decayTop}\n|A|_s^2 = \\sum_{j \\in \\mathbb Z, l \\in \\mathbb Z^\\nu} \\big( \\sup_{j_1 - j_2 = j} |A_{j_1}^{j_2}(l)| \\big)^2 \\langle l,j \\rangle^{2 s} \\, . \n\\end{equation}\nThese matrices are identified with the $ \\varphi $-dependent family\nof operators\n\\begin{equation}\\label{Aphi}\nA(\\varphi) := \\big( A_{j_1}^{j_2} (\\varphi)\\big)_{j_1, j_2 \\in \\mathbb Z} \\, , \\quad \nA_{j_1}^{j_2} (\\varphi) := {\\mathop\\sum}_{l \\in \\mathbb Z^\\nu} A_{j_1}^{j_2}(l) e^{{\\mathrm i} l \\cdot \\varphi}\n\\end{equation}\nwhich act on functions of the $x$-variable as\n\\begin{equation}\\label{notationA}\nA(\\varphi) : h(x) = \\sum_{j \\in \\mathbb Z} h_j e^{{\\mathrm i} jx} \\mapsto \nA(\\varphi) h(x) = \\sum_{j_1, j_2 \\in \\mathbb Z} A_{j_1}^{j_2} (\\varphi) h_{j_2} e^{{\\mathrm i} j_1 x} \\, . \n\\end{equation}\nWe still denote by $ | A(\\varphi) |_s $ the $ s $-decay norm of the matrix in \\eqref{Aphi}.\nAs in \\cite{BBM}, all the transformations that we shall construct in this paper are of this type (with $ j, j_1, j_2 \\neq 0 $ because\nthey act on the phase space $ H^1_0 (\\mathbb T_x) $).\nThis observation allows to interpret the conjugacy procedure from a dynamical point of view, see \\cite{BBM}-section 2.2. \nLet us fix some terminology.\n\n\\begin{definition}\\label{operatore Hamiltoniano}\nWe say that:\n\nthe operator \n$(A h)(\\varphi, x) := A(\\varphi) h(\\varphi, x)$ is {\\sc symplectic} if each \n$ A (\\varphi ) $, $ \\varphi \\in \\mathbb T^\\nu $, is a symplectic map of the phase space\n(or of a symplectic subspace like $ H_S^\\bot $);\n\n\nthe operator\n$\n\\omega \\cdot \\partial_{\\varphi} - \\partial_x G( \\varphi ) $ is {\\sc Hamiltonian} \nif each $ G (\\varphi) $, $ \\varphi \\in \\mathbb T^\\nu $, is symmetric; \n\nan operator is {\\sc real} if it maps real-valued functions into real-valued functions.\n\\end{definition}\n\n\nAs well known, a Hamiltonian operator $ \\omega \\cdot \\partial_{\\varphi} - \\partial_x G( \\varphi ) $ is transformed,\nunder a symplectic map $ {\\cal A} $, \ninto another Hamiltonian operator $ \\omega \\cdot \\partial_{\\varphi} - \\partial_x E( \\varphi ) $, \nsee e.g. \\cite{BBM}-section 2.3.\n\n\\smallskip\n\nWe conclude this preliminary section recalling the following well known lemmata, see Appendix of \\cite{BBM}.\n\n\\begin{lemma} {\\bf (Composition)}\n\\label{lemma:composition of functions, Moser}\nAssume $ f \\in C^s (\\mathbb T^d \\times B_1)$, $B_1 := \\{ y \\in \\mathbb R^m : \n|y| \\leq 1 \\}$. Then \n$ \\forall u \\in H^{s}(\\mathbb T^d, \\mathbb R^m) $ such that $ \\| u \\|_{L^\\infty} < 1 $, \nthe composition operator $\\tilde{f}(u)(x) := f(x, u(x))$ satisfies\n$ \\| \\tilde f(u) \\|_s \\leq C \\| f \\|_{C^s} (\\|u\\|_{s} + 1) $\nwhere the constant $C $ depends on $ s ,d $. If $ f \\in C^{s+2} $ \nand $ \\| u + h \\|_{L^\\infty} < 1$, then \n\\begin{align*}\n\\big\\| \\tilde f(u+h) - {\\mathop\\sum}_{i = 0}^k \\frac{\\tilde{f}^{(i)}(u)}{i !} [h^i] \\big\\|_s\n& \\leq C \\| f \\|_{C^{s+ 2}} \\, \\| h \\|_{L^\\infty}^k ( \\| h \\|_{s} + \\| h \\|_{L^\\infty} \\| u \\|_{s}) \\, , \\quad k = 0, 1 \\, . \n\\end{align*}\nThe previous statement also holds replacing \n$\\| \\ \\|_s$ \nwith the norms $| \\ |_{s, \\infty} $. \n\\end{lemma}\n\n\\begin{lemma} \\label{lemma:standard Sobolev norms properties}\n{\\bf (Tame product).} \nFor $s \\geq s_0 > d\/2 $, \n$$\n\\| uv \\|_s \\leq C(s_0) \\|u\\|_s \\|v\\|_{s_0} + C(s) \\|u\\|_{s_0} \\| v \\|_s \\, , \n\\quad \\forall u,v \\in H^s(\\mathbb T^d) \\, .\n$$\nFor $s \\geq 0$, $s \\in \\mathbb N$,\n$$\n\\| uv \\|_s \\leq \\tfrac32 \\, \\| u \\|_ {L^\\infty} \\| v \\|_s + C(s) \\| u \\|_ {W^{s, \\infty}} \\| v \\|_0 \\, , \n\\quad \\forall u \\in W^{s,\\infty}(\\mathbb T^d) \\, , \\ v \\in H^s(\\mathbb T^d) \\, .\n$$\nThe above inequalities also hold for \nthe norms $\\Vert \\ \\Vert_s^{{\\rm Lip}(\\gamma)}$.\n\\end{lemma} \n\n\n\\begin{lemma} {\\bf (Change of variable)} \\label{lemma:utile} \nLet $p \\in W^{s,\\infty} (\\mathbb T^d,\\mathbb R^d) $, $ s \\geq 1$, with \n$ \\| p \\|_{W^{1, \\infty}} \\leq 1\/2 $. Then the function $f(x) = x + p(x)$\nis invertible, with inverse $ f^{-1}(y) = y + q(y)$ where\n $q \\in W^{s,\\infty}(\\mathbb T^d,\\mathbb R^d)$, and \n $ \\| q \\|_{W^{s, \\infty}} \\leq C \\| p \\|_{ W^{s, \\infty}} $.\nIf, moreover, $p = p_\\omega $ depends in a Lipschitz way on \na parameter $\\omega \\in \\Omega \\subset \\mathbb R^\\nu $, \nand \n$ \\| D_x p_\\omega \\|_ {L^\\infty} \\leq 1\/2 $, $ \\forall \\omega $, \nthen\n$ \\| q \\|_{W^{s, \\infty}}^{{\\rm Lip}(\\gamma)} \\leq C \\| p \\|_{W^{s+1, \\infty}}^{{\\rm Lip}(\\gamma)} $.\nThe constant $C := C (d, s) $ is independent of $\\gamma$.\n\nIf $u \\in H^s (\\mathbb T^d,\\mathbb C)$, then $ (u\\circ f)(x) := u(x+p(x))$ satisfies \n\\begin{align*}\n\\| u \\circ f \\|_s \n& \\leq C (\\|u\\|_s + \\| p \\|_{W^{s, \\infty}} \\|u\\|_1),\n\\quad\n\\| u \\circ f - u \\|_s \n\\leq C ( \\| p \\|_{L^\\infty} \\| u \\|_{s + 1} + \\| p \\|_{W^{s, \\infty}} \\| u \\|_{2} ) ,\n\\\\\n\\| u \\circ f \\|_{s}^{{{\\rm Lip}(\\gamma)}}\n& \\leq C \\, \n\\big( \\| u \\|_{s+1}^{{{\\rm Lip}(\\gamma)}} + \\| p \\|_{W^{s, \\infty}}^{{\\rm Lip}(\\gamma)}\\| u \\|_2^{{\\rm Lip}(\\gamma)} \\big). \\nonumber\n\\end{align*}\nThe function $u \\circ f^{-1} $ satisfies the same bounds.\n\\end{lemma}\n\n\n\n\\section{Weak Birkhoff normal form}\\label{sec:WBNF}\n\nThe Hamiltonian\nof the perturbed KdV equation \\eqref{kdv quadratica} is $ H = H_2 + H_3 + H_{\\geq 5} $ (see \\eqref{Ham in intro}) where \n\\begin{equation} \\label{H iniziale KdV}\nH_2 (u):=\\frac{1}{2} \\int_{\\mathbb T} u_x^{2} \\, dx \\, , \\quad \nH_3(u) := \\int_\\mathbb T u^3 dx \\, , \\quad \nH_{\\geq 5}(u) := \\int_\\mathbb T f(x, u,u_x) dx \\,,\n\\end{equation}\nand $f$ satisfies \\eqref{order5}. \nAccording to the splitting \\eqref{u = v + z} $ u = v + z $, $ v \\in H_S $, $ z \\in H_S^\\bot $, \nwe have\n\\begin{equation}\\label{prima v z}\nH_2(u) = \\int_\\mathbb T \\frac{v_x^2}{2}\\, dx + \\int_\\mathbb T \\frac{z_x^2}{2} \\, dx, \n\\quad \nH_3 (u) = \\int_{\\mathbb T} v^3 dx + 3 \\int_{\\mathbb T} v^2 z dx \n+ 3\\int_{\\mathbb T} v z^2 dx + \\int_{\\mathbb T} z^3 dx \\, . \n\\end{equation}\nFor a finite-dimensional space\n\\begin{equation} \\label{def E finito}\nE := E_{C} := \\mathrm{span} \\{ e^{{\\mathrm i} jx} : 0 < |j| \\leq C \\}, \\quad C > 0,\n\\end{equation}\nlet $\\Pi_E $ denote the corresponding $ L^2 $-projector on $E$.\n\nThe notation $R(v^{k-q} z^q)$ indicates a homogeneous polynomial of degree $k$ in $(v,z)$ of the form \n$$\nR(v^{k-q} z^q) = M[\\underbrace{v, \\ldots, v}_{(k-q)\\, \\text{times}}, \n\\underbrace{z, \\ldots, z}_{q \\, \\text{times}} \\,], \\qquad \nM = k\\text{-linear} \\, .\n$$\n\\begin{proposition} \\label{prop:weak BNF} \n{\\bf (Weak Birkhoff normal form)} Assume Hypothesis $ ({\\mathtt S}2) $. \nThen there exists an analytic invertible symplectic transformation \nof the phase space $ \\Phi_B : H^1_0 (\\mathbb T_x) \\to H^1_0 (\\mathbb T_x) $ \nof the form \n\\begin{equation} \\label{finito finito}\n\\Phi_B(u) = u + \\Psi(u), \n\\quad\n\\Psi(u) = \\Pi_E \\Psi(\\Pi_E u),\n\\end{equation}\nwhere $ E $ is a finite-dimensional space as in \\eqref{def E finito}, such that the transformed Hamiltonian is\n\\begin{equation} \\label{widetilde cal H}\n{\\cal H} := H \\circ \\Phi_B \n= H_2 + \\mathcal{H}_3 + \\mathcal{H}_4 + {\\cal H}_5 + {\\cal H}_{\\geq 6} \\,,\n\\end{equation}\nwhere $H_2$ is defined in \\eqref{H iniziale KdV}, \n\\begin{equation} \\label{H3tilde} \n\\mathcal{H}_3 := \\int_{\\mathbb T} z^3\\,dx + 3 \\int_{\\mathbb T} v z^2 \\,dx \\,, \\quad \n\\mathcal{H}_4 := - \\frac32 \\sum_{j \\in S} \\frac{|u_j|^4}{j^2} + \\mathcal{H}_{4,2} + \\mathcal{H}_{4,3} \\,, \\quad \n{\\cal H}_5 := \\sum_{q=2}^5 R(v^{5-q} z^q)\\,,\n\\end{equation} \n\\begin{equation} \\label{mH3 mH4}\n\\mathcal{H}_{4,2} := 6 \\int_\\mathbb T v z \\Pi_S \\big((\\partial_x^{-1} v)(\\partial_x^{-1} z) \\big)\\,dx + 3 \\int_\\mathbb T z^2 \\pi_0 (\\partial_x^{-1} v)^2 \\,dx\\,, \\quad \n\\mathcal{H}_{4, 3} := R(v z^3) \\,, \n\\end{equation}\nand ${\\cal H}_{\\geq 6}$ collects all the terms of order at least six in $(v,z)$. \n\\end{proposition}\n\nThe rest of this section is devoted to the proof of Proposition \\ref{prop:weak BNF}.\n\n\n\nFirst, we remove the cubic terms $ \\int_{\\mathbb T} v^3 + 3 \\int_{\\mathbb T} v^2 z $ from the Hamiltonian $ H_3 $ defined in \\eqref{prima v z}.\nIn the Fourier coordinates \\eqref{Fourier}, we have\n\\begin{equation}\\label{H3 Fourier}\nH_2 = \\frac12 \\sum_{j \\neq 0} j^2 |u_j|^2, \n\\quad H_3 = \\sum_{j_1 + j_2 + j_3 = 0} u_{j_1} u_{j_2} u_{j_3} \\, .\n\\end{equation}\nWe look for a symplectic transformation $ \\Phi^{(3)} $ of the phase space \nwhich eliminates the monomials $ u_{j_1} u_{j_2} u_{j_3} $ of $ H_3 $ with at most {\\it one} index\noutside $ S $. Note that, by the relation $ j_1 + j_2 + j_3 = 0 $, they are {\\it finitely} many.\nWe look for $\\Phi^{(3)} := (\\Phi^t_{F^{(3)}})_{|t=1}$ as the time-1 flow map generated by the Hamiltonian vector field $X_{F^{(3)}}$, with an auxiliary Hamiltonian of the form\n$$\nF^{(3)}(u) := \\sum_{j_1 + j_2 + j_3 = 0} F^{(3)}_{j_1 j_2 j_3} u_{j_1} u_{j_2} u_{j_3} \\,.\n$$\nThe transformed Hamiltonian is \n\\begin{align}\nH^{(3)} & := H \\circ \\Phi^{(3)} \n= H_2 + H_3^{(3)} + H_4^{(3)} + H_{\\geq 5}^{(3)} \\,,\n\\notag \\\\ \n\\label{H tilde 234}\nH_3^{(3)} & = H_3 + \\{ H_2, F^{(3)} \\}, \\quad \n H_4^{(3)} = \\frac12 \\{ \\{ H_2, F^{(3)}\\}, F^{(3)}\\} + \\{H_3, F^{(3)} \\} , \n\\end{align}\nwhere $ H_{\\geq 5}^{(3)} $ collects all the terms of order at least five in $(u,u_x)$.\nBy \\eqref{H3 Fourier} and \\eqref{PoissonBr} we calculate \n$$\nH_3^{(3)} = \n\\sum_{j_1 + j_2 + j_3 = 0} \n\\big\\{ 1 - {\\mathrm i} (j_1^3 + j_2^3 + j_3^3) F^{(3)}_{j_1 j_2 j_3} \\big\\} \\, u_{j_1} u_{j_2} u_{j_3} \\,.\n$$\nHence, in order to eliminate the monomials with at most one index outside $ S $, \nwe choose\n\\begin{equation}\\label{F3q}\nF^{(3)}_{j_1 j_2 j_3} := \\begin{cases}\n \\dfrac{1}{{\\mathrm i} (j_1^3 + j_2^3 + j_3^3)} & \\text{if} \\,\\,(j_1,j_2,j_3) \\in {\\cal A}\\,, \n\\\\\n0 & \\text{otherwise},\n\\end{cases}\n\\end{equation}\nwhere ${\\cal A} := \\big\\{ (j_1 , j_2 , j_3) \\in (\\mathbb Z \\setminus \\{ 0 \\})^3$ : $j_1 + j_2 + j_3 = 0$, \n$j_1^{3} + j_2^{3} + j_3^{3} \\neq 0$, and at least 2 among $j_1 , j_2 , j_3$ belong to $S \\big\\}$. \nNote that \n\\begin{equation} \\label{def calA quadratica}\n{\\cal A} = \\big\\{ (j_1 , j_2 , j_3) \\in (\\mathbb Z \\setminus \\{ 0 \\})^3 : \nj_1 + j_2 + j_3 = 0, \\, \\text{and at least 2 among}\\,\\, j_1 , j_2 , j_3 \\,\\, \n\\text{belong to} \\, S \\big\\} \n\\end{equation}\nbecause of the elementary relation\n\\begin{equation}\\label{prodottino}\nj_1 + j_2 + j_3 = 0 \\quad \\Rightarrow \\quad j_1^3 + j_2^3 + j_3^3 = 3 j_1 j_2 j_3 \\neq 0 \n\\end{equation}\nbeing $ j_1, j_2, j_3 \\in \\mathbb Z \\setminus \\{ 0 \\}$. \nAlso note that $ \\mathcal{A} $ is a finite set, actually $ \\mathcal{A} \\subseteq [- 2 C_S, 2 C_S]^{3} $ where \nthe tangential sites $ S \\subseteq [- C_S, C_S ]$.\nAs a consequence, the Hamiltonian vector field $ X_{F^{(3)}} $ has finite rank and vanishes outside \nthe finite dimensional subspace $ E := E_{2 C_S}$ (see \\eqref{def E finito}), namely\n$$\n X_{ F^{(3)}}(u) = \\Pi_E X_{ F^{(3)}} ( \\Pi_E u ) \\, . \n$$\nHence its flow $ \\Phi^{(3)} : H^1_0 (\\mathbb T_x) \\to H^1_0 (\\mathbb T_x) $ has the form \\eqref{finito finito} and it is analytic. \n\nBy construction, all the monomials of $ H_3 $ \nwith at least two indices outside $ S $ \nare not modified by the transformation $ \\Phi^{(3)} $. Hence (see \\eqref{prima v z}) we have \n\\begin{equation}\\label{lem:H3tilde}\nH_3^{(3)} = \\int_{\\mathbb T} z^3\\,dx + 3 \\int_{\\mathbb T} v z^2 \\,dx \\, .\n\\end{equation}\nWe now compute the fourth order term\n$ H_4^{(3)} = \\sum_{i=0}^4 H_{4,i}^{(3)} $ in \\eqref{H tilde 234}, \nwhere $ H_{4,i}^{(3)} $ is of type $ R( v^{4-i} z^i )$.\n\n\\begin{lemma}\nOne has (recall the definition \\eqref{def pi 0} of $ \\pi_0 $) \n\\begin{equation}\\label{Htilde41} \n{H}_{4,0} ^{(3)} := \\frac32 \\int_\\mathbb T v^2 \\pi_0[(\\partial_x^{-1} v)^2] dx \\, , \\quad\nH_{4,2}^{(3)} := 6 \\int_\\mathbb T v z \\Pi_S \\big((\\partial_x^{-1} v)(\\partial_x^{-1} z) \\big)\\,dx \n+ 3 \\int_\\mathbb T z^2 \\pi_0 [(\\partial_x^{-1} v)^2] dx \\, . \n\\end{equation}\n\\end{lemma}\n\n\\begin{proof}\nWe write $ H_3 = H_{3, \\leq 1} + H_3^{(3)} $ where $ H_{3, \\leq 1}(u) := \\int_\\mathbb T v^3 dx + 3 \\int_\\mathbb T v^2 z \\, dx $. \nThen, by \\eqref{H tilde 234}, we get\n\\begin{equation} \\label{grado 4 *}\nH_4^{(3)} \n= \\frac12 \t\\big\\{ H_{3, \\leq 1} \\, , F^{(3)} \\big\\} + \\{ H_3^{(3)} , F^{(3)} \\} \\,.\n\\end{equation} \nBy \\eqref{F3q}, \\eqref{prodottino}, the auxiliary Hamiltonian may be written as\n$$\nF^{(3)} (u) \n= - \\frac{1}{3}\\sum_{(j_1, j_2, j_3) \\in {\\cal A}} \n\\frac{u_{j_1} u_{j_2} u_{j_3}}{ ({\\mathrm i} j_1) ( {\\mathrm i} j_2) ( {\\mathrm i} j_3)} \n= - \\frac{1}{3}\\int_\\mathbb T (\\partial_x^{-1} v)^3 dx - \\int_\\mathbb T (\\partial_x^{-1} v)^2 (\\partial_x^{-1} z) dx \\, . \n$$\nHence, using that the projectors $\\Pi_S$, $\\Pi_S^\\bot $ are self-adjoint and $\\partial_x^{-1}$ is skew-selfadjoint, \n\\begin{equation} \\label{grad F3 formula}\n\\nabla F^{(3)}(u) \n= \\partial_x^{-1}\\big\\{ (\\partial_x^{-1} v)^2 + 2 \\Pi_S \\big[ (\\partial_x^{-1} v)(\\partial_x^{-1} z)\\big] \\big\\}\n\\end{equation}\n(we have used that $\\partial_x^{-1} \\pi_0 = \\partial_x^{-1}$ be the definition of $\\partial_x^{-1}$). \nRecalling the Poisson bracket definition \\eqref{Poisson bracket}, \nusing that $ \\nabla H_{3, \\leq 1}(u) = 3 v^2 + 6 \\Pi_S(v z) $ and \\eqref{grad F3 formula},\nwe get \n\\begin{align}\n\\{ H_{3, \\leq1}, F^{(3)} \\} \n& = \\int_\\mathbb T \\big\\{ 3 v^2 + 6 \\Pi_S(v z) \\big\\} \\pi_0 \\big\\{ (\\partial_x^{-1} v)^2 + \n2 \\Pi_S \\big[ (\\partial_x^{-1} v)(\\partial_x^{-1} z)\\big] \\big\\}\\,dx \n\\nonumber\n\\\\\n& = 3 \\int_\\mathbb T v^2 \\pi_0 (\\partial_x^{-1} v)^2\\,dx +\n 12 \\int_\\mathbb T \\Pi_S(v z) \\Pi_S [ (\\partial_x^{-1} v)(\\partial_x^{-1} z) ]\\,dx + R(v^3 z) \\, . \\label{H31F}\n\\end{align}\nSimilarly, since $ \\nabla H_3^{(3)}(u) = 3 z^2 + 6 \\Pi_S^\\bot (v z) $, \n\\begin{equation}\n\\{ H_3^{(3)}, F^{(3)} \\} \n= \n3 \\int_\\mathbb T z^2 \\pi_0 (\\partial_x^{-1} v)^2\\,dx \n+ R(v^3 z) + R(v z^3) \\,. \\label{H3tildeF}\n\\end{equation}\nThe lemma follows by \\eqref{grado 4 *}, \\eqref{H31F}, \\eqref{H3tildeF}.\n\\end{proof}\n\nWe now construct a symplectic map $ \\Phi^{(4)} $ \nsuch that the Hamiltonian system obtained transforming $ H_2 + H_3^{(3)} + H_4^{(3)} $ \npossesses the invariant subspace $ H_S $ (see \\eqref{splitting S-S-bot}) \nand its dynamics on $ H_S $ is integrable and non-isocronous. \nHence we have to eliminate the term $ H_{4,1}^{(3)} $ \n(which is linear in $ z $),\nand to normalize $ H_{4,0}^{(3)} $ (which is independent of $ z $). \nWe need the following elementary lemma (Lemma 13.4 in \\cite{KaP}).\n\n\\begin{lemma} \\label{lemma:interi} \nLet $j_1, j_2, j_3, j_4 \\in \\mathbb Z $ such that $ j_1 + j_2 + j_3 + j_4 = 0 $. \nThen \n$$\nj_1^3 + j_2^3 + j_3^3 + j_4^3 = -3 (j_1 + j_2) (j_1 + j_3) (j_2 + j_3).\n$$\n\\end{lemma}\n\n\n\\begin{lemma}\nThere exists a symplectic transformation $ \\Phi^{(4)} $ of the form \\eqref{finito finito}\nsuch that\n\\begin{equation}\\label{Ham4quadratica}\nH^{(4)} := H^{(3)} \\circ \\Phi^{(4)} = \nH_2 + H_3^{(3)} + H_4^{(4)} + H^{(4)}_{\\geq 5} \\,, \n\\qquad \nH^{(4)}_4 := - \\frac32 \\sum_{j \\in S} \\frac{|u_j|^4}{j^2} + H_{4,2}^{(3)} + H_{4,3}^{(3)} \\,,\n\\end{equation}\nwhere $ H_3^{(3)} $ is defined in \\eqref{lem:H3tilde},\n$ H_{4,2}^{(3)} $ in \\eqref{Htilde41}, $ H_{4,3}^{(3)} = R( v z^3) $ \nand $ H_{\\geq 5}^{(4)} $ collects all the terms of degree at least five in $(u,u_x)$.\n\\end{lemma}\n\n\\begin{proof}\nWe look for a map $\\Phi^{(4)} := (\\Phi_{F^{(4)}}^t)_{|t=1}$ which is the time $ 1$-flow map \nof an auxiliary Hamiltonian\n$$\nF^{(4)}(u) := \\sum_{\\begin{subarray}{c}\nj_1 + j_2 + j_3 + j_4 = 0 \\\\\n\\text{at least}\\,\\,3\\,\\,\\text{indices are in}\\,\\,S\n\\end{subarray}} F^{(4)}_{j_1 j_2 j_3 j_4} u_{j_1} u_{j_2} u_{j_3} u_{j_4} \n$$\nwith the same form of the Hamiltonian $ H_{4,0}^{(3)} + H_{4,1}^{(3)} $.\nThe transformed Hamiltonian is\n\\begin{equation}\\label{callH}\nH^{(4)} := H^{(3)} \\circ \\Phi^{(4)} \n= H_2 + H_3^{(3)} + H_4^{(4)} + H_{\\geq 5}^{(4)}, \\quad \nH_4^{(4)} = \\{ H_2, F^{(4)} \\} + H_4^{(3)} , \n\\end{equation}\nwhere $ H_{\\geq 5}^{(4)} $ collects all the terms of order at least five. \nWe write $ H_4^{(4)} = \\sum_{i=0}^4 H_{4,i}^{(4)} $ where each $ H_{4,i}^{(4)} $ if of type $ R( v^{4-i} z^i ) $. \nWe choose the coefficients \n\\begin{equation}\\label{def F4}\nF^{(4)}_{j_1 j_2 j_3 j_4} := \\begin{cases}\n \\dfrac{ H^{(3)}_{j_1j_2j_3j_4} }{{\\mathrm i} (j_1^3 + j_2^3 + j_3^3+ j_4^3)} & \\text{if} \\,\\,(j_1,j_2,j_3, j_4) \\in {\\cal A}_4\\,, \n\\\\\n0 & \\text{otherwise},\n\\end{cases}\n\\end{equation}\nwhere \n\\begin{align*} \n{\\cal A}_4 := \\big\\{ (j_1 , j_2 , j_3, j_4) \\in (\\mathbb Z \\setminus \\{ 0 \\})^4 & : \nj_1 + j_2 + j_3 + j_4 = 0, \\, j_1^{3} + j_2^{3} + j_3^{3} + j_4^{3} \\neq 0, \\\\\n& \\quad \\text{and at most one among}\\,\\, j_1 , j_2 , j_3, j_4 \\,\\, \\text{outside} \\, S \\big\\} \\, . \n\\end{align*}\nBy this definition $ H_{4,1}^{(4)}= 0 $\nbecause there exist no integers $ j_1, j_2 , j_3 \\in S$, $ j_4 \\in S^c $ satisfying \n$ j_1 + j_2 + j_3 + j_4 = 0 $, $ j_1^3 + j_2^3 + j_3^3 + j_4^3 = 0 $, by Lemma \\ref{lemma:interi} and the\nfact that $ S $ is symmetric. \nBy construction, the terms $ H_{4,i}^{(4)}= H_{4,i}^{(3)} $, $ i = 2, 3, 4$, are not changed by $ \\Phi^{(4)} $.\nFinally, by \\eqref{Htilde41} \n\\begin{equation} \\label{H2F1}\nH_{4,0}^{(4)}= \\frac32 \n\\sum_{\\begin{subarray}{c}\nj_1, j_2, j_3, j_4 \\in S \\\\\nj_1 + j_2 + j_3 + j_4 = 0 \\\\\nj_1^3 + j_2^3 + j_3^3 + j_4^3 = 0 \\\\\nj_1 + j_2 \\,,\\,j_3 + j_4 \\neq 0\n\\end{subarray}} \\frac{1}{({\\mathrm i} j_3) ({\\mathrm i} j_4)}u_{j_1} u_{j_2} u_{j_3} u_{j_4} \\, . \n\\end{equation}\nIf $ j_1 + j_2 + j_3 + j_4 = 0 $ and $ j_1^3 + j_2^3 + j_3^3 + j_4^3 = 0$, then $(j_1 + j_2)(j_1 + j_3)(j_2 + j_3) = 0$ by Lemma \\ref{lemma:interi}. \nWe develop the sum in \\eqref{H2F1} with respect to the first index $j_1$. Since $ j_1 + j_2 \\neq 0 $\nthe possible cases are:\n\\[\n(i) \\ \\big\\{ \nj_2 \\neq - j_1, \\ \nj_3 = - j_1, \\\nj_4 = - j_2 \\big\\} \n\\qquad \\text{or} \\qquad\n(ii) \\ \\big\\{\nj_2 \\neq - j_1, \\ \nj_3 \\neq - j_1, \\ \nj_3 = - j_2, \\ \nj_4 = - j_1 \\big\\} . \n\\]\nHence, using $ u_{-j} = \\bar{u}_j $ (recall \\eqref{Fourier}), and since $S$ is symmetric, we have\n\\begin{equation} \\label{case ii}\n\\sum_{(i)} \\frac{1}{j_3 j_4} u_{j_1} u_{j_2} u_{j_3} u_{j_4} \n= \\sum_{j_1, j_2 \\in S, j_2 \\neq - j_1} \\frac{|u_{j_1}|^2 |u_{j_2}|^2 }{j_1 j_2} \n= \\sum_{j,j' \\in S} \\frac{|u_j|^2 |u_{j'}|^2}{j j'} + \\sum_{j \\in S} \\frac{|u_j|^4}{j^2} = \\sum_{j \\in S} \\frac{|u_j|^4}{j^2}\\,, \n\\end{equation}\nand in the second case ($ii$) \n\\begin{equation} \\label{case iii}\n\\sum_{(ii)} \\frac{1}{j_3 j_4} u_{j_1} u_{j_2} u_{j_3} u_{j_4} \n = \\sum_{j_1, j_2, j_2 \\neq \\pm j_1} \\frac{1}{j_1 j_2} u_{j_1} u_{j_2} u_{-j_2} u_{-j_1} =\n \\sum_{j \\in S} \\frac{1}{j} |u_{j}|^2 \\Big( \\sum_{j_2 \\neq \\pm j} \\frac{1}{j_2} |u_{j_2}|^2 \\Big) = 0\\,. \n\\end{equation}\nThen \\eqref{Ham4quadratica} follows by \\eqref{H2F1}, \\eqref{case ii}, \\eqref{case iii}. \n\\end{proof}\n\nNote that the Hamiltonian $ H_2 + H_3^{(3)} + H_4^{(4)} $ (see \\eqref{Ham4quadratica}) possesses\nthe invariant subspace $ \\{ z = 0 \\} $ and the system restricted to $ \\{ z = 0 \\} $ is\ncompletely integrable and non-isochronous (actually it is formed by $ \\nu $ decoupled rotators). \nWe shall construct quasi-periodic solutions which bifurcate from this invariant manifold. \n\nIn order to enter in a perturbative regime, \nwe have to eliminate further monomials of $ H^{(4)} $ in \\eqref{Ham4quadratica}. \nThe minimal requirement for the convergence of the nonlinear Nash-Moser iteration is to \neliminate the monomials $ R(v^5) $ and $ R(v^4 z)$.\nHere we need the choice of the sites of Hypothesis $ ({\\mathtt S}2) $. \n\n\\begin{remark}\\label{remark:cubic}\nIn the KAM theorems \\cite{k1}, \\cite{Po3} (and \\cite{PP}, \\cite{Wang}),\nas well as for the perturbed mKdV equations \\eqref{mKdV}, \nthese further steps of Birkhoff normal form are not required because \nthe nonlinearity of the original PDE is yet cubic. A difficulty of KdV\nis that the nonlinearity is quadratic. \n\\end{remark}\n\nWe spell out Hypothesis $( {\\mathtt S}2)$ as follows:\n\\begin{itemize}\n\\item {\\sc $( {\\mathtt S}2_0)$.} \nThere is no choice of $ 5 $ integers $ j_1, \\ldots, j_5 \\in S $ such that \n\\begin{equation}\\label{scelta dei siti grado 7}\nj_1 + \\ldots + j_5 = 0\\,,\\quad \nj_1^3 + \\ldots + j_5^3 = 0 \\,.\n\\end{equation} \n\\item {\\sc $({\\mathtt S}2_1)$.} \nThere is no choice of $4 $ integers $j_1, \\ldots, j_{4} $ in $ S $ \nand an integer in the complementary set $ j_5 \\in S^c := (\\mathbb Z \\setminus \\{0\\}) \\setminus S $ such that \\eqref{scelta dei siti grado 7} holds. \n\\end{itemize}\n\nThe homogeneous component of degree $ 5 $ of $H^{(4)}$ is\n$$\nH^{(4)}_5 (u) = \\sum_{j_1+ \\ldots + j_5 = 0} H^{(4)}_{j_1, \\ldots, j_5} u_{j_1} \\ldots u_{j_5} \\,.\n$$\nWe want to remove from $H^{(4)}_5$ the terms with at most one index among $ j_1, \\ldots , j_5 $ outside $ S $.\nWe consider the auxiliary Hamiltonian \n\\begin{equation} \\label{Fn odd}\nF^{(5)} = \\sum_{\\begin{subarray}{c} j_1 + \\ldots + j_5 = 0 \\\\ \n\\text{at most one index outside $ S $}\n\\end{subarray} } \nF_{j_1 \\ldots j_5}^{(5)} u_{j_1} \\ldots u_{j_5} \\, , \\quad \n F_{j_1 \\ldots j_5}^{(5)} := \\frac{H_{j_1 \\ldots j_5}^{(5)}}{ {\\mathrm i} (j_1^3 + \\ldots + j_5^3)} \\,.\n\\end{equation}\nBy Hypotheses $ ({\\mathtt S}2_0), ({\\mathtt S}2_1) $, if $ j_1 + \\ldots + j_5 = 0 $ with at most one index outside $ S $ then \n$ j_1^3 + \\ldots + j_5^3 \\neq 0 $ and \n$ F^{(5)}$ is well defined. \nLet $ \\Phi^{(5)} $ be the time $ 1 $-flow generated by \n$ X_{F^{(5)}} $. The new Hamiltonian is\n\\begin{equation}\\label{callHn}\nH^{(5)} := H^{(4)} \\circ \\Phi^{(5)} \n= H_2 + H_3^{(3)} + H_4^{(4)} + \\{ H_2, F^{(5)} \\} + H_5^{(4)} + H^{(5)}_{\\geq 6} \n\\end{equation}\nwhere, by \\eqref{Fn odd}, \n$$\nH_5^{(5)} := \\{ H_2, F^{(5)} \\} + H_5^{(4)} = {\\mathop\\sum}_{q = 2}^5 R(v^{5 - q} z^q) \\, .\n$$\n\n\nRenaming $ {\\cal H} := H^{(5)} $, namely $ {\\cal H}_n := H^{(n)}_n $, $ n =3, 4, 5 $, and \nsetting $ \\Phi_B := \\Phi^{(3)} \\circ \\Phi^{(4)} \\circ \\Phi^{(5)} $, formula \\eqref{widetilde cal H} follows. \n\nThe homogeneous component $ H^{(4)}_{5} $ preserves the momentum, see section \\ref{sec:Ham For}. \nHence $F^{(5)} $ also preserves the momentum. \nAs a consequence, also $ H^{(5)}_k $, $ k \\leq 5 $, preserve the momentum. \n\n\nFinally, since $F^{(5)} $ is Fourier-supported on a finite set, \nthe transformation $\\Phi^{(5)}$ is of type \\eqref{finito finito} (and analytic), \nand therefore also the composition $ \\Phi_B $ is of type \\eqref{finito finito} (and analytic). \n\n\\section{Action-angle variables}\\label{sec:4}\n\n\nWe now introduce action-angle variables on the tangential directions by the change of coordinates\n\\begin{equation}\\label{coordinate azione angolo}\n\\begin{cases}\nu_j := \\sqrt{\\xi_j + |j| y_j} \\, e^{{\\mathrm i} \\theta_j}, \\qquad & \\text{if} \\ j \\in S\\,,\\\\\nu_j := z_j, \\qquad & \\text{if} \\ j \\in S^c \\, , \n\\end{cases}\n\\end{equation}\nwhere (recall $ u_{-j} = {\\overline u}_j $)\n\\begin{equation}\\label{simmeS}\n\\xi_{-j} = \\xi_j \\, , \\quad \n\\xi_j > 0 \\, , \\quad \ny_{-j} = y_j \\, , \\quad \n\\theta_{-j} = - \\theta_j \\, , \\quad \n\\theta_j, \\, y_j \\in \\mathbb R \\, , \\quad \n\\forall j \\in S \\,.\n\\end{equation}\nFor the tangential sites $ S^+ := \\{ {\\bar \\jmath_1}, \\ldots, {\\bar \\jmath_\\nu} \\} $ we shall also denote \n$ \\theta_{\\bar \\jmath_i} := \\theta_i $, $ y_{\\bar \\jmath_i} := y_i $, \n$ \\xi_{\\bar \\jmath_i} := \\xi_i $, $ i =1, \\ldots \\, \\nu $.\n\nThe symplectic 2-form $ \\Omega $ in \\eqref{2form0} (i.e. \\eqref{2form KdV}) becomes \n\\begin{equation}\\label{2form}\n{\\cal W} := \\sum_{i=1}^\\nu d \\theta_i \\wedge d y_i \n+ \\frac12 \\sum_{j \\in S^c \\setminus \\{ 0 \\} } \\frac{1}{{\\mathrm i} j} \\, d z_j \\wedge d z_{-j} = \n\\big( \\sum_{i=1}^\\nu d \\theta_i \\wedge d y_i \\big) \\oplus \\Omega_{S^\\bot} = d \\Lambda \n\\end{equation}\nwhere $ \\Omega_{S^\\bot} $ denotes the restriction of $ \\Omega $ to $ H_S^\\bot $ (see \\eqref{splitting S-S-bot}) and \n$ \\Lambda $ is the contact $ 1 $-form on $ \\mathbb T^\\nu \\times \\mathbb R^\\nu \\times H_S^\\bot $ defined by\n$ \\Lambda_{(\\theta, y, z)} : \\mathbb R^\\nu \\times \\mathbb R^\\nu \\times H_S^\\bot \\to \\mathbb R $, \n\\begin{equation}\\label{Lambda 1 form}\n\\Lambda_{(\\theta, y, z)}[\\widehat \\theta, \\widehat y, \\widehat z] := \n- y \\cdot \\widehat \\theta + \\frac12 ( \\partial_x^{-1} z, \\widehat z )_{L^2 (\\mathbb T)} \\, . \n\\end{equation}\nInstead of working in a shrinking neighborhood of the origin, it is a convenient devise to rescale the ``unperturbed actions\" $ \\xi $\nand the action-angle variables as\n\\begin{equation}\\label{rescaling kdv quadratica}\n \\xi \\mapsto \\varepsilon^2 \\xi \\, , \\quad y \\mapsto \\varepsilon^{2b} y \\, , \\quad z \\mapsto \\varepsilon^b z \\, . \n \\end{equation}\nThen the symplectic $ 2 $-form in \\eqref{2form} transforms into $ \\varepsilon^{2b} {\\cal W } $. Hence\nthe Hamiltonian system generated by $ {\\cal H} $ in \\eqref{widetilde cal H} \ntransforms into the new Hamiltonian system \n\\begin{equation} \\label{def H eps}\n \\dot \\theta = \\partial_y H_\\varepsilon (\\theta, y, z) \\, , \\ \n \\dot y = - \\partial_\\theta H_\\varepsilon (\\theta, y, z) \\, , \\ \n z_t = \\partial_x \\nabla_z H_\\varepsilon (\\theta, y, z) \\, , \\quad H_\\varepsilon := \\varepsilon^{-2b} \\mathcal{H} \\circ A_\\varepsilon\n\\end{equation}\nwhere \n\\begin{equation} \\label{def A eps}\nA_\\varepsilon (\\theta, y, z) := \\varepsilon v_\\varepsilon(\\theta, y) + \\varepsilon^b z \n:= \\varepsilon {\\mathop \\sum}_{j \\in S} \\sqrt{\\xi_j + \\varepsilon^{2(b-1)} |j| y_j} \\, e^{{\\mathrm i} \\theta_j} e^{{\\mathrm i} j x} + \\varepsilon^b z \\, .\n\\end{equation}\nWe shall still denote by $ X_{H_\\varepsilon} = (\\partial_y H_\\varepsilon, - \\partial_\\theta H_\\varepsilon, \\partial_x \\nabla_z H_\\varepsilon) $ the Hamiltonian vector field in the variables $ (\\theta, y, z ) \\in \\mathbb T^\\nu \\times \\mathbb R^\\nu \\times H_S^\\bot $.\n\n\nWe now write explicitly the Hamiltonian $ H_\\varepsilon (\\theta, y, z) $ in \\eqref{def H eps}. \nThe quadratic Hamiltonian $ H_2 $ in \\eqref{H iniziale KdV} transforms into\n\\begin{equation}\\label{shape H2}\n\\varepsilon^{-2b}H_2 \\circ A_\\varepsilon = const + \n{\\mathop \\sum}_{j \\in S^+} j^3 y_j + \\frac{1}{2} \\int_{\\mathbb T} z_x^2 \\, dx \\, , \n\\end{equation}\nand, recalling \\eqref{H3tilde}, \\eqref{mH3 mH4}, \nthe Hamiltonian $ {\\cal H }$ in \\eqref{widetilde cal H} \ntransforms into\n(shortly writing $ v_\\varepsilon := v_\\varepsilon (\\theta, y) $) \n\\begin{align}\nH_\\varepsilon (\\theta, y, z) & \n= e(\\xi) + \\a (\\xi) \\cdot y + \\frac12 \\int_\\mathbb T z_x^2 dx + \\varepsilon^b\\int_\\mathbb T z^3 dx + 3 \\varepsilon \\int_\\mathbb T v_\\varepsilon z^2 dx \\label{formaHep}\n\\\\ & \\quad \n+ \\varepsilon^2 \\Big\\{6 \\int_\\mathbb T v_\\varepsilon z\n \\Pi_S \\big((\\partial_x^{-1} v_\\varepsilon)(\\partial_x^{-1} z) \\big)\\,dx + 3 \\int_\\mathbb T z^2 \\pi_0 (\\partial_x^{-1} v_\\varepsilon)^2\\,dx \\Big\\} \n - \\frac32 \\varepsilon^{2b} {\\mathop \\sum}_{j \\in S} y_j^2 \\nonumber \n\\\\ & \\quad\n+ \\varepsilon^{b + 1} R(v_\\varepsilon z^3) \n+ \\varepsilon^3 R(v_\\varepsilon^3 z^2) \n+ \\varepsilon^{2+b} \\sum_{q=3}^5 \\varepsilon^{(q-3)(b-1)} R(v_\\varepsilon^{5-q} z^q) \n+ \\varepsilon^{-2b} {\\cal H}_{\\geq 6} (\\varepsilon v_\\varepsilon + \\varepsilon^b z ) \\nonumber \n\\end{align}\nwhere $e(\\xi)$ is a constant, and the frequency-amplitude map is\n\\begin{equation} \\label{mappa freq amp}\n\\a(\\xi) := \\bar\\omega + \\varepsilon^2 {\\mathbb A} \\xi \\, , \\quad \n{\\mathbb A} := - 6 \\, \\text{diag} \\{ 1\/j \\}_{j \\in S^+} \\, .\n\\end{equation} \nWe write the Hamiltonian in \\eqref{formaHep} as \n\\begin{equation} \\label{Hamiltoniana Heps KdV}\nH_\\varepsilon = {\\cal N} + P \\,, \\quad \n{\\cal N}(\\theta, y, z) \n= \\a (\\xi) \\cdot y + \\frac12 \\big(N(\\theta) z , z \\big)_{L^2(\\mathbb T)} \\,,\n\\end{equation}\nwhere \n\\begin{align}\n\\frac12 \\big(N(\\theta) z, z \\big)_{L^2(\\mathbb T)} & \n:= \\frac12 \\big( (\\partial_{z} \\nabla H_\\varepsilon)(\\theta, 0, 0)[z], z \\big)_{L^2(\\mathbb T)} \\label{Nshape} \n= \\frac12 \\int_\\mathbb T z_x^2 dx + 3 \\varepsilon \\int_\\mathbb T v_\\varepsilon(\\theta, 0) z^2 dx \\\\\n& + \\varepsilon^2 \\Big\\{6 \\int_\\mathbb T v_\\varepsilon(\\theta, 0) z \\Pi_S \n\\big((\\partial_x^{-1} v_\\varepsilon(\\theta,0))(\\partial_x^{-1} z) \\big) dx \n+ 3 \\int_\\mathbb T z^2 \\pi_0 (\\partial_x^{-1} v_\\varepsilon(\\theta, 0))^2 dx \\Big\\} + \\ldots \\nonumber\n\\end{align}\nand $P := H_\\varepsilon - {\\cal N} $.\n\n\\section{The nonlinear functional setting}\\label{sec:functional}\n\nWe look for an embedded invariant torus \n\\begin{equation} \\label{embedded torus i}\ni : \\mathbb T^\\nu \\to \\mathbb T^\\nu \\times \\mathbb R^\\nu \\times H_S^\\bot, \\quad \n\\varphi \\mapsto i (\\varphi) := ( \\theta (\\varphi), y (\\varphi), z (\\varphi)) \n\\end{equation}\nof the Hamiltonian vector field $ X_{H_\\varepsilon} $ \nfilled by quasi-periodic solutions with diophantine frequency $ \\omega $. We require that $ \\omega $ belongs to the set \n\\begin{equation}\\label{Omega epsilon}\n\\Omega_\\varepsilon := \\a ( [1,2]^\\nu ) = \n\\{ \\a(\\xi) : \\xi \\in [1,2]^\\nu \\} \n\\end{equation}\nwhere $ \\a $ is the diffeomorphism \n\\eqref{mappa freq amp}, and, in the Hamiltonian $ H_\\varepsilon $ in \\eqref{Hamiltoniana Heps KdV}, we choose \n\\begin{equation}\\label{linkxiomega}\n\\xi = \\a^{-1}(\\omega) = \\varepsilon^{-2} {\\mathbb A}^{-1} (\\omega - \\bar\\omega) \\, . \n\\end{equation}\nSince any $ \\omega \\in \\Omega_\\varepsilon$ is $ \\varepsilon^2 $-close to the integer vector\n$ \\bar \\omega $ (see \\eqref{bar omega}), we require that the constant $\\gamma$ in the diophantine inequality \n\\begin{equation}\\label{omdio}\n |\\omega \\cdot l | \\geq \\gamma \\langle l \\rangle^{-\\tau} \\, , \\ \\ \\forall l \\in \\mathbb Z^\\nu \\setminus \\{0\\} \\, , \n\\quad \\text{satisfies} \\ \\ \\gamma = \\varepsilon^{2+a} \\quad \\text{for some} \\ a > 0 \\,. \n\\end{equation}\nWe remark that the definition of $\\gamma$ in \\eqref{omdio} is slightly stronger than the minimal condition, which is $ \\gamma \\leq c \\varepsilon^2 $ with $ c $ small enough.\nIn addition to \\eqref{omdio} we shall also require that $ \\omega $ satisfies the first and second order Melnikov-non-resonance conditions \\eqref{Omegainfty}. \n\nWe\nlook for an embedded invariant torus \nof the modified Hamiltonian vector field $ X_{H_{\\varepsilon, \\zeta}} = X_{H_\\varepsilon} + (0, \\zeta, 0)$ which is generated by \nthe Hamiltonian \n\\begin{equation}\\label{hamiltoniana modificata}\nH_{\\varepsilon, \\zeta} (\\theta, y, z) := H_\\varepsilon (\\theta, y, z) + \\zeta \\cdot \\theta\\,,\\quad \\zeta \\in \\mathbb R^\\nu\\,.\n\\end{equation}\nNote that $ X_{H_{\\varepsilon, \\zeta}} $ is periodic in $\\theta $\n(unlike $ H_{\\varepsilon, \\zeta} $). \nIt turns out that an invariant torus for $ X_{H_{\\varepsilon, \\zeta}} $ is actually \ninvariant for $ X_{H_\\varepsilon} $, see Lemma \\ref{zeta = 0}. \nWe introduce the parameter $\\zeta \\in \\mathbb R^\\nu $ in order to control the average in the \n$ y$-component\nof the linearized equations.\nThus we look for zeros of the nonlinear operator\n\\begin{align}\n {\\cal F} (i, \\zeta ) \n& := {\\cal F} (i, \\zeta, \\omega, \\varepsilon ) := {\\cal D}_\\omega i (\\varphi) - X_{H_{\\varepsilon, \\zeta}} (i(\\varphi)) =\n{\\cal D}_\\omega i (\\varphi) - X_{\\cal N} (i(\\varphi)) - X_P (i(\\varphi)) + (0, \\zeta, 0 ) \\label{operatorF} \\\\\n& \\nonumber := \\left(\n\\begin{array}{c}\n{\\cal D}_\\omega \\theta (\\varphi) - \\partial_y H_\\varepsilon ( i(\\varphi) ) \\\\\n{\\cal D}_\\omega y (\\varphi) + \\partial_\\theta H_\\varepsilon ( i(\\varphi) ) + \\zeta \\\\\n{\\cal D}_\\omega z (\\varphi) - \\partial_x \\nabla_z H_\\varepsilon ( i(\\varphi)) \n\\end{array}\n\\right) \\!\\! \n= \\!\\!\n\\left(\n\\begin{array}{c} \n{\\cal D}_\\omega \\Theta (\\varphi) - \\partial_y P (i(\\varphi) ) \\\\\n\\!\\! \\! {\\cal D}_\\omega y (\\varphi) + \\frac12 \\partial_\\theta ( N(\\theta (\\varphi)) z(\\varphi), z(\\varphi) )_{L^2(\\mathbb T)} \n\t+ \\partial_\\theta P ( i(\\varphi) ) + \\zeta \\!\\!\\!\\! \\\\\n{\\cal D}_\\omega z (\\varphi) - \\partial_x N ( \\theta (\\varphi )) z (\\varphi) \n- \\partial_x \\nabla_z P ( i(\\varphi) ) \n\\end{array} \n\\right) \n\\end{align}\nwhere $ \\Theta(\\varphi) := \\theta (\\varphi) - \\varphi $ is $ (2 \\pi)^\\nu $-periodic \nand we use the short notation \n\\begin{equation}\\label{Domega}\n{\\cal D}_\\omega := \\omega \\cdot \\partial_\\varphi \\, . \n\\end{equation}\nThe Sobolev norm of the periodic component of the embedded torus \n\\begin{equation}\\label{componente periodica}\n{\\mathfrak I}(\\varphi) := i (\\varphi) - (\\varphi,0,0) := ( {\\Theta} (\\varphi), y(\\varphi), z(\\varphi))\\,, \\quad \\Theta(\\varphi) := \\theta (\\varphi) - \\varphi \\, , \n\\end{equation}\nis \n\\begin{equation}\\label{norma fracchia}\n\\| {\\mathfrak I} \\|_s := \\| \\Theta \\|_{H^s_\\varphi} + \\| y \\|_{H^s_\\varphi} + \\| z \\|_s \n\\end{equation}\nwhere $ \\| z \\|_s := \\| z \\|_{H^s_{\\varphi,x}} $ is defined in \\eqref{Sobolev coppia}.\nWe link the rescaling \\eqref{rescaling kdv quadratica}\nwith the diophantine constant $ \\gamma = \\varepsilon^{2+a} $ by choosing\n\\begin{equation}\\label{link gamma b}\n\\gamma = \\varepsilon^{2b}\\,, \\qquad\nb = 1 + ( a \\slash 2 ) \\, .\n\\end{equation}\nOther choices are possible, \nsee Remark \\ref{comm3}.\n\n\\begin{theorem}\\label{main theorem}\nLet the tangential sites $ S $ in \\eqref{tang sites} \nsatisfy\n$({\\mathtt S}1), ({\\mathtt S}2)$. \nThen, for all $ \\varepsilon \\in (0, \\varepsilon_0 ) $, where $ \\varepsilon_0 $ is small enough, \nthere exists a Cantor-like set $ {\\cal C}_\\varepsilon \\subset \\Omega_\\varepsilon $,\nwith asympotically full measure as $ \\varepsilon \\to 0 $, namely\n\\begin{equation}\\label{stima in misura main theorem}\n\\lim_{\\varepsilon\\to 0} \\, \\frac{|{\\cal C}_\\varepsilon|}{|\\Omega_\\varepsilon|} = 1 \\, , \n\\end{equation}\nsuch that, for all $ \\omega \\in {\\cal C}_\\varepsilon $, there exists a solution $ i_\\infty (\\varphi) := i_\\infty (\\omega, \\varepsilon)(\\varphi) $ \nof ${\\cal D}_\\omega i_\\infty(\\varphi) - X_{H_\\varepsilon}(i_\\infty(\\varphi)) = 0 $.\nHence the embedded torus \n$ \\varphi \\mapsto i_\\infty (\\varphi) $ is invariant for the Hamiltonian vector field $ X_{H_\\varepsilon (\\cdot, \\xi)} $ \nwith $ \\xi $ as in \\eqref{linkxiomega}, and it is filled by quasi-periodic solutions with frequency $ \\omega $.\nThe torus $i_\\infty$ satisfies\n\\begin{equation}\\label{stima toro finale}\n\\| i_\\infty (\\varphi) - (\\varphi,0,0) \\|_{s_0 + \\mu}^{\\mathrm{Lip}(\\g)} = O(\\varepsilon^{6 - 2 b} \\gamma^{-1} ) \n\\end{equation}\nfor some $ \\mu := \\mu (\\nu) > 0 $. Moreover, the torus $ i_\\infty $ is {\\sc linearly stable}. \n\\end{theorem}\n\n\nTheorem \\ref{main theorem} is proved in sections \\ref{costruzione dell'inverso approssimato}-\\ref{sec:NM}. It implies Theorem \\ref{thm:KdV}\nwhere the $ \\xi_j $ in \\eqref{solution u} are $ \\varepsilon^2 \\xi_j $, $ \\xi_j \\in [1,2 ] $, in \\eqref{linkxiomega}.\nBy \\eqref{stima toro finale}, \ngoing back to the variables before the rescaling \\eqref{rescaling kdv quadratica}, \nwe get\n$ \\Theta_\\infty = O( \\varepsilon^{6-2b} \\gamma^{-1}) $, $ y_\\infty = O( \\varepsilon^6 \\gamma^{-1} ) $, $ z_\\infty = O( \\varepsilon^{6-b} \\gamma^{-1} ) $, which, \nas $ b \\to 1^+ $, tend to the expected optimal estimates. \n\n\\begin{remark} \\label{comm3}\nThere are other possible ways to link the rescaling \\eqref{rescaling kdv quadratica}\nwith the diophantine constant $ \\gamma = \\varepsilon^{2+a} $.\nThe choice $ \\gamma > \\varepsilon^{2b} $\nreduces to study perturbations of an isochronous system (as in \\cite{Ku}, \\cite{k1}, \\cite{Po3}), and it is convenient to \n introduce $ \\xi (\\omega) $ as a variable. \nThe case $ \\varepsilon^{2b} > \\gamma $, in particular $ b = 1 $, \nhas to be dealt with a perturbation approach of a non-isochronous system {\\`a la} Arnold-Kolmogorov. \n\\end{remark}\n\n\nWe now give the tame estimates for the composition operator induced by the Hamiltonian vector fields $ X_{\\cal N} $ and $ X_P $ in \\eqref{operatorF}, that we shall use in the next sections. \n\nWe first estimate the composition operator induced by $ v_\\varepsilon (\\theta, y) $ defined in \\eqref{def A eps}.\nSince the functions $ y \\mapsto \\sqrt{\\xi + \\varepsilon^{2(b - 1)}|j| y} $, $\\theta \\mapsto e^{{\\mathrm i} \\theta}$ \nare analytic for $\\varepsilon$ small enough and $|y| \\leq C$, \nthe composition Lemma \\ref{lemma:composition of functions, Moser} implies that, for all $ \\Theta, y \\in H^s(\\mathbb T^\\nu, \\mathbb R^\\nu )$, \n$ \\| \\Theta \\|_{s_0}, \\| y \\|_{s_0} \\leq 1 $, setting $\\theta(\\varphi) := \\varphi + \\Theta (\\varphi)$, \n$ \\| v_\\varepsilon (\\theta (\\varphi) ,y(\\varphi) ) \\|_s \n\\leq_s 1 + \\| \\Theta \\|_s + \\| y \\|_s $. \nHence, using also \n\\eqref{linkxiomega}, the map $ A_\\varepsilon $ in \\eqref{def A eps} satisfies, for all \n$ \\| {\\mathfrak I} \\|_{s_0}^{\\mathrm{Lip}(\\g)} \\leq 1 $ (see \\eqref{componente periodica})\n\\begin{equation} \\label{stima Aep}\n \\| A_\\varepsilon (\\theta (\\varphi),y(\\varphi),z(\\varphi)) \\|_s^{{\\mathrm{Lip}(\\g)}} \\leq_s \\varepsilon (1 + \\| {\\mathfrak I} \\|_s^{\\mathrm{Lip}(\\g)}) \\, . \n\\end{equation} \nWe now give tame estimates for the Hamiltonian vector fields $ X_{\\cal N} $, $ X_P $, $ X_{H_\\varepsilon} $, \nsee \\eqref{Hamiltoniana Heps KdV}-\\eqref{Nshape}.\n\n\\begin{lemma}\\label{lemma quantitativo forma normale}\nLet $ {\\mathfrak{I}}(\\varphi) $ in \\eqref{componente periodica} satisfy\n$ \\| {\\mathfrak I} \\|_{s_0 + 3}^{\\mathrm{Lip}(\\g)} \\leq C\\varepsilon^{6 - 2b} \\gamma^{-1} $. \nThen \n\\begin{alignat}{2} \n\\| \\partial_y P(i) \\|_s^{\\mathrm{Lip}(\\g)} & \\leq_s \\varepsilon^4 + \\varepsilon^{2b} \\| {\\mathfrak I}\\|_{s+1}^{\\mathrm{Lip}(\\g)} \\label{D y P} \\, , \n& \n\\| \\partial_\\theta P(i) \\|_s^{\\mathrm{Lip}(\\g)}\n& \\leq_s \\varepsilon^{6 - 2b} (1 + \\| {\\mathfrak I} \\|_{s + 1}^{\\mathrm{Lip}(\\g)}) \n\\\\ \n\\| \\nabla_z P(i) \\|_s^{\\mathrm{Lip}(\\g)}\n& \\leq_s \\varepsilon^{5 - b} + \\varepsilon^{6- b} \\gamma^{-1} \\| {\\mathfrak I} \\|_{s + 1}^{\\mathrm{Lip}(\\g)} \\, ,\n&\n\\| X_P(i)\\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^{6 - 2b} + \\varepsilon^{2b} \\| {\\mathfrak I}\\|_{s + 3}^{\\mathrm{Lip}(\\g)} \\label{stima XP}\n\\\\\n\\label{stime linearizzato campo hamiltoniano}\n\\| \\partial_{\\theta} \\partial_y P(i)\\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^{4} + \\varepsilon^{5} \\gamma^{-1} \\| {\\mathfrak I}\\|_{s + 2}^{\\mathrm{Lip}(\\g)} \\, ,\n& \n\\| \\partial_y \\nabla_z P(i)\\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^{b+3} + \\varepsilon^{2b - 1} \\| {\\mathfrak I} \\|_{s+2}^{\\mathrm{Lip}(\\g)} \n\\end{alignat}\n\\begin{equation}\n\\label{D yy P}\n\\| \\partial_{yy} P(i) + 3 \\varepsilon^{2b} I_{\\mathbb R^\\nu} \\|_s^{\\mathrm{Lip}(\\g)} \n\\leq_s \\varepsilon^{2+ 2b} + \\varepsilon^{2b+3} \\gamma^{-1} \\| {\\mathfrak{I}} \\|_{s+2}^{\\mathrm{Lip}(\\g)}\n\\end{equation}\nand, for all $ \\widehat \\imath := (\\widehat \\Theta, \\widehat y, \\widehat z) $, \n\\begin{align}\n\\label{D yii}\n\\| \\partial_y d_{i} X_P(i)[\\widehat \\imath \\,]\\|_s^{\\mathrm{Lip}(\\g)} &\\leq_s \\varepsilon^{2 b - 1} \\big( \\| \\widehat \\imath \\|_{s + 3}^{\\mathrm{Lip}(\\g)} + \\| {\\mathfrak I}\\|_{s + 3}^{\\mathrm{Lip}(\\g)} \\| \\widehat \\imath \\|_{s_0 + 3}^{\\mathrm{Lip}(\\g)}\\big) \\\\\n\\label{tame commutatori} \\| d_i X_{H_\\varepsilon}(i) [\\widehat \\imath \\, ] + (0,0, \\partial_{xxx} \\hat z)\\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon \\big( \\| \\widehat \\imath\\|_{s + 3}^{\\mathrm{Lip}(\\g)} \n+ \\|{\\mathfrak I} \\|_{s + 3}^{\\mathrm{Lip}(\\g)} \\| \\widehat \\imath\\|_{s_0 + 3}^{\\mathrm{Lip}(\\g)} \\big) \\\\\n\\label{parte quadratica da P}\n\\| d_i^2 X_{H_\\varepsilon}(i) [\\widehat \\imath, \\widehat \\imath \\,]\\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon\\Big( \\| \\widehat \\imath \\|_{s + 3}^{\\mathrm{Lip}(\\g)} \\| \n\\widehat \\imath\\|_{s_0 + 3}^{\\mathrm{Lip}(\\g)} + \\| {\\mathfrak I}\\|_{s + 3}^{\\mathrm{Lip}(\\g)} (\\| \\widehat \\imath\\|_{s_0 + 3}^{\\mathrm{Lip}(\\g)})^2\\Big).\n\\end{align}\n\\end{lemma}\n\nIn the sequel we will also use that, by the diophantine condition \\eqref{omdio}, the operator $ {\\cal D}_\\omega^{-1} $ (see \\eqref{Domega})\nis defined for all functions $ u $ with zero $ \\varphi $-average, and satisfies \n\\begin{equation}\\label{Dom inverso}\n \\| {\\cal D}_\\omega^{-1} u \\|_s \\leq C \\gamma^{-1} \\| u \\|_{s+ \\tau} \\, , \\quad \\| {\\cal D}_\\omega^{-1} u \\|_s^{{\\mathrm{Lip}(\\g)}} \\leq C \\gamma^{-1} \\| u \\|_{s+ 2 \\tau+1}^{{\\mathrm{Lip}(\\g)}} \\, . \n\\end{equation}\n\n\\section{Approximate inverse}\\label{costruzione dell'inverso approssimato}\n\nIn order to implement a convergent Nash-Moser scheme that leads to a solution of \n$ \\mathcal{F}(i, \\zeta) = 0 $ \n our aim is to construct an \\emph{approximate right inverse} (which satisfies tame estimates) of the linearized operator \n\\begin{equation}\\label{operatore linearizzato}\nd_{i, \\zeta} {\\cal F}(i_0, \\zeta_0 )[\\widehat \\imath \\,, \\widehat \\zeta ] =\nd_{i, \\zeta} {\\cal F}(i_0 )[\\widehat \\imath \\,, \\widehat \\zeta ] = \n{\\cal D}_\\omega \\widehat \\imath - d_i X_{H_\\varepsilon} ( i_0 (\\varphi) ) [\\widehat \\imath ] + (0, \\widehat \\zeta, 0 )\n \\,, \n\\end{equation}\nsee Theorem \\ref{thm:stima inverso approssimato}. Note that \n$ d_{i, \\zeta} {\\cal F}(i_0, \\zeta_0 ) = d_{i, \\zeta} {\\cal F}(i_0 ) $ is independent of $ \\zeta_0 $ (see \\eqref{operatorF}).\n\nThe notion of approximate right inverse is introduced in \\cite{Z1}. It denotes a linear operator \nwhich is an \\emph{exact} right inverse at a solution $ (i_0, \\zeta_0) $ of $ {\\cal F}(i_0, \\zeta_0) = 0 $.\nWe want to implement the general strategy in \\cite{BB13}-\\cite{BB14}\nwhich reduces the search of an approximate right inverse of \\eqref{operatore linearizzato} \nto the search of an approximate inverse on the normal directions only. \n\n\nIt is well known that an invariant torus $ i_0 $ with diophantine flow \nis isotropic (see e.g. \\cite{BB13}), namely the pull-back $ 1$-form $ i_0^* \\Lambda $ is closed, \nwhere $ \\Lambda $ is the contact 1-form in \\eqref{Lambda 1 form}. \nThis is tantamount to say that the 2-form $ \\cal W $ (see \\eqref{2form}) vanishes on the torus $ i_0 (\\mathbb T^\\nu )$ \n(i.e. $\\cal W$ vanishes on the tangent space at each point $i_0(\\varphi)$ of the manifold \n$i_0(\\mathbb T^\\nu)$), because \n$ i_0^* {\\cal W} = i_0^* d \\Lambda = d i_0^* \\Lambda $. \nFor an ``approximately invariant\" torus $ i_0 $ the 1-form $ i_0^* \\Lambda$ is only ``approximately closed\".\nIn order to make this statement quantitative we consider\n\\begin{equation}\\label{coefficienti pull back di Lambda}\ni_0^* \\Lambda = {\\mathop \\sum}_{k = 1}^\\nu a_k (\\varphi) d \\varphi_k \\,,\\quad \na_k(\\varphi) := - \\big( [\\partial_\\varphi \\theta_0 (\\varphi)]^T y_0 (\\varphi) \\big)_k \n+ \\frac12 ( \\partial_{\\varphi_k} z_0(\\varphi), \\partial_{x}^{-1} z_0(\\varphi) )_{L^2(\\mathbb T)}\n\\end{equation}\nand we quantify how small is \n\\begin{equation} \\label{def Akj} \ni_0^* {\\cal W} = d \\, i_0^* \\Lambda = {\\mathop\\sum}_{1 \\leq k < j \\leq \\nu} A_{k j}(\\varphi) d \\varphi_k \\wedge d \\varphi_j\\,,\\quad A_{k j} (\\varphi) := \n\\partial_{\\varphi_k} a_j(\\varphi) - \\partial_{\\varphi_j} a_k(\\varphi) \\, .\n\\end{equation}\nAlong this section we will always assume the following hypothesis \n(which will be verified at each step of the Nash-Moser iteration):\n\n\\begin{itemize}\n\\item {\\sc Assumption.} \nThe map $\\omega\\mapsto i_0(\\omega)$ is a Lipschitz function defined on some subset $\\Omega_o \\subset \\Omega_\\varepsilon$, where $\\Omega_\\varepsilon$ is defined in \\eqref{Omega epsilon}, and, for some $ \\mu := \\mu (\\t, \\nu) > 0 $, \n\\begin{equation}\\label{ansatz 0}\n\\| {\\mathfrak I}_0 \\|_{s_0+\\mu}^{{\\mathrm{Lip}(\\g)}} \\leq C\\varepsilon^{6 - 2b} \\gamma^{-1}, \\quad \n\\| Z \\|_{s_0 + \\mu}^{{\\mathrm{Lip}(\\g)}} \\leq C \\varepsilon^{6 - 2b}, \\quad \n\\gamma = \\varepsilon^{2 + a}, \n\\quad \n b := 1 + (a\/2) \\,,\n \\quad a \\in (0, 1 \/ 6), \n\\end{equation}\nwhere ${\\mathfrak{I}}_0(\\varphi) := i_0(\\varphi) - (\\varphi,0,0)$, and \n\\begin{equation} \\label{def Zetone}\nZ(\\varphi) := (Z_1, Z_2, Z_3) (\\varphi) := {\\cal F}(i_0, \\zeta_0) (\\varphi) =\n\\omega \\cdot \\partial_\\varphi i_0(\\varphi) - X_{H_{\\varepsilon, \\zeta_0}}(i_0(\\varphi)) \\, .\n\\end{equation}\n\\end{itemize}\n\\begin{lemma} \\label{zeta = 0}\n$ |\\zeta_0|^{{\\mathrm{Lip}(\\g)}} \\leq C \\| Z \\|_{s_0}^{{\\mathrm{Lip}(\\g)}}$ \\!\\!\\!. If $ {\\cal F}(i_0, \\zeta_0) = 0 $ then $ \\zeta_0 = 0 $, \n namely the torus $i_0 $ is invariant for $X_{H_\\varepsilon}$.\n \\end{lemma}\n\n\\begin{proof}\nIt is proved in \\cite{BB13} the formula\n$$\n\\zeta_0 = \n\\int_{\\mathbb T^\\nu} - [\\partial_\\varphi y_0 (\\varphi)]^T Z_1 (\\varphi) \n+ [\\partial_\\varphi \\theta_0 (\\varphi)]^T Z_2 (\\varphi) \n- [\\partial_\\varphi z_0 (\\varphi) ]^T \\partial_x^{-1} Z_3 (\\varphi) \\, d \\varphi \\, .\n$$\nHence the lemma follows by \\eqref{ansatz 0} and usual algebra estimate. \n\\end{proof}\n\nWe now quantify the size of $ i_0^* {\\cal W} $ in terms of $ Z $.\n\n\\begin{lemma}\nThe coefficients $A_{kj} (\\varphi) $ in \\eqref{def Akj} satisfy \n\\begin{equation}\\label{stima A ij}\n\\| A_{k j} \\|_s^{{\\mathrm{Lip}(\\g)}} \\leq_s \\gamma^{-1} \\big(\\| Z \\|_{s+2\\t+2}^{{\\mathrm{Lip}(\\g)}} \\| {\\mathfrak I}_0 \\|_{s_0+ 1}^{{\\mathrm{Lip}(\\g)}} \n+ \\| Z \\|_{s_0+1}^{{\\mathrm{Lip}(\\g)}} \\| {\\mathfrak I}_0 \\|_{s+ 2 \\tau + 2}^{{\\mathrm{Lip}(\\g)}} \\big)\\,.\n\\end{equation}\n\\end{lemma}\n\n\\begin{proof}\nWe estimate the coefficients of the Lie derivative\n$ L_\\omega (i_0^* {\\cal W}) := \\sum_{k < j} {\\cal D}_\\omega A_{k j}(\\varphi) d \\varphi_k \\wedge d \\varphi_j $. \nDenoting by $ \\underline{e}_k $ the $ k $-th versor of $ \\mathbb R^\\nu $ we have \n$$\n {\\cal D}_\\omega A_{k j} \n= L_\\omega( i_0^* {\\cal W})(\\varphi)[ \\underline{e}_k , \\underline{e}_j ] \n= {\\cal W}\\big( \\partial_\\varphi Z(\\varphi) \\underline{e}_k , \\partial_\\varphi i_0(\\varphi) \\underline{e}_j \\big) \n+ {\\cal W} \\big(\\partial_\\varphi i_0(\\varphi) \\underline{e}_k , \\partial_\\varphi Z(\\varphi) \\underline{e}_j \\big) \n$$\n(see \\cite{BB13}). \nHence\n\\begin{equation} \\label{bella trovata}\n\\| {\\cal D}_\\omega A_{k j} \\|_s^{\\mathrm{Lip}(\\g)} \n\\leq_s \\| Z \\|_{s+1}^{\\mathrm{Lip}(\\g)} \\| {\\mathfrak I}_0 \\|_{s_0 + 1}^{\\mathrm{Lip}(\\g)} \n+ \\| Z \\|_{s_0 + 1}^{\\mathrm{Lip}(\\g)} \\| {\\mathfrak I}_0 \\|_{s + 1}^{\\mathrm{Lip}(\\g)} \\,. \n\\end{equation}\nThe bound \\eqref{stima A ij} follows applying $ {\\cal D}_\\omega^{-1}$ and using \n\\eqref{def Akj}, \\eqref{Dom inverso}. \n\\end{proof}\n\nAs in \\cite{BB13} we first modify the approximate torus $ i_0 $ to obtain an isotropic torus $ i_\\d $ which is \nstill approximately invariant. We denote the Laplacian $ \\Delta_\\varphi := \\sum_{k=1}^\\nu \\partial_{\\varphi_k}^2 $ . \n\n\\begin{lemma}\\label{toro isotropico modificato} {\\bf (Isotropic torus)} \nThe torus $ i_\\delta(\\varphi) := (\\theta_0(\\varphi), y_\\delta(\\varphi), z_0(\\varphi) ) $ defined by \n\\begin{equation}\\label{y 0 - y delta}\ny_\\d := y_0 + [\\partial_\\varphi \\theta_0(\\varphi)]^{- T} \\rho(\\varphi) \\, , \\qquad \n\\rho_j(\\varphi) := \\Delta_\\varphi^{-1} {\\mathop\\sum}_{ k = 1}^\\nu \\partial_{\\varphi_j} A_{k j}(\\varphi) \n\\end{equation}\nis isotropic. \nIf \\eqref{ansatz 0} holds, then, for some $ \\sigma := \\sigma(\\nu,\\t) $, \n\\begin{align} \\label{stima y - y delta}\n\\| y_\\delta - y_0 \\|_s^{{\\mathrm{Lip}(\\g)}} \n& \\leq_s \\gamma^{-1} \\big(\\| Z \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\| {\\mathfrak I}_0 \\|_{s_0 + \\sigma}^{{\\mathrm{Lip}(\\g)}} + \n\\| Z \\|_{s_0 + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\| {\\mathfrak I}_0 \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\big) \\,,\n\\\\\n\\label{stima toro modificato}\n\\| {\\cal F}(i_\\delta, \\zeta_0) \\|_s^{{\\mathrm{Lip}(\\g)}} \n& \\leq_s \\| Z \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} + \\| Z \\|_{s_0 + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\| {\\mathfrak I}_0 \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\\\\n\\label{derivata i delta}\n\\| \\partial_i [ i_\\d][ \\widehat \\imath ] \\|_s & \\leq_s \\| \\widehat \\imath \\|_s + \\| {\\mathfrak I}_0\\|_{s + \\sigma} \\| \\widehat \\imath \\|_s \\, .\n\\end{align}\n\\end{lemma}\nIn the paper we denote equivalently the differential by $ \\partial_i $ or $ d_i $. Moreover we denote \nby $ \\sigma := \\sigma(\\nu, \\tau ) $ possibly different (larger) ``loss of derivatives\" constants. \n\n\\begin{proof}\nIn this proof we write $\\| \\ \\|_s$ to denote $\\| \\ \\|_s^{{\\mathrm{Lip}(\\g)}}$.\nThe proof of the isotropy of $i_\\d$ is in \\cite{BB13}. \nThe estimate \\eqref{stima y - y delta} follows by \\eqref{y 0 - y delta}, \\eqref{stima A ij}, \\eqref{ansatz 0}\nand the tame bound for the inverse \n$ \\Vert [\\partial_\\varphi \\theta_0 ]^{-T}\\Vert_{s} \\leq_s 1 + \\| {\\mathfrak I}_0 \\|_{s + 1} $. \nIt remains to estimate the difference (see \\eqref{operatorF} and note that $ X_{\\cal N} $ does not depend on $ y $)\n\\begin{equation}\\label{F diff}\n{\\cal F}(i_\\delta, \\zeta_0) - {\\cal F}(i_0, \\zeta_0) = \n\\begin{pmatrix} 0 \\\\ {\\cal D}_\\omega (y_\\delta - y_0 ) \\\\ 0 \\end{pmatrix}\\, \n+ X_P(i_\\delta ) - X_P(i_0).\n\\end{equation}\nUsing \\eqref{stime linearizzato campo hamiltoniano}, \\eqref{D yy P}, we get \n$ \\| \\partial_y X_P (i) \\|_s \\leq_s \\varepsilon^{2b} + \\varepsilon^{2b - 1} \\| {\\mathfrak I}\\|_{s + 3} $.\nHence \\eqref{stima y - y delta}, \n\\eqref{ansatz 0} imply\n\\begin{equation}\\label{XP delta - XP}\n\\|X_ P(i_\\delta ) - X_P(i_0 )\\|_s \n\\leq_s \\| {\\mathfrak I}_0 \\|_{s_0 + \\sigma} \\|Z \\|_{s + \\sigma} + \\| {\\mathfrak I}_0 \\|_{s + \\sigma} \\|Z \\|_{s_0 + \\sigma} \\, . \n\\end{equation}\nDifferentiating \\eqref{y 0 - y delta} we have \n\\begin{equation}\\label{D omega y 0 - y delta}\n{\\cal D}_\\omega (y_\\delta - y_0 ) \n= [\\partial_\\varphi \\theta_0(\\varphi)]^{-T} {\\cal D}_\\omega \\rho(\\varphi) \n+ ( {\\cal D}_\\omega [\\partial_\\varphi \\theta_0(\\varphi)]^{-T} ) \\rho(\\varphi) \n\\end{equation}\nand \n$ {\\cal D}_\\omega \\rho_j(\\varphi) = \\Delta^{-1}_\\varphi \\sum_{k = 1}^\\nu \\partial_{\\varphi_j} {\\cal D}_\\omega A_{kj}(\\varphi) $. \nUsing \\eqref{bella trovata}, we deduce that \n\\begin{equation}\\label{primo pezzo D omega y 0 - y delta}\n\\| [\\partial_\\varphi \\theta_0 ]^{-T} {\\cal D}_\\omega \\rho \\|_s \n\\leq_s \\| Z\\|_{s + 1} \\| {\\mathfrak I}_0\\|_{s_0 + 1} + \\| Z \\|_{s_0 + 1} \\| {\\mathfrak I}_0\\|_{s + 1}\\,.\n\\end{equation}\nTo estimate the second term in \\eqref{D omega y 0 - y delta}, \nwe differentiate \n$ Z_1(\\varphi) = {\\cal D}_\\omega \\theta_0(\\varphi) - \\omega - (\\partial_y P)(i_0(\\varphi)) $\n(which is the first component in \\eqref{operatorF}) with respect to $\\varphi$. We get\n$ {\\cal D}_\\omega \\partial_\\varphi \\theta_0(\\varphi) \n= \\partial_\\varphi (\\partial_y P)(i_0(\\varphi)) + \\partial_\\varphi Z_1(\\varphi) $. \nThen, by \\eqref{D y P},\n\\begin{equation}\\label{bella trovata 2}\n\\|{\\cal D}_\\omega [\\partial_\\varphi \\theta_0]^T \\|_s \n\\leq_s \\varepsilon^4 + \\varepsilon^{2b} \\| {\\mathfrak I}_0 \\|_{s + 2} + \\| Z \\|_{s + 1}\\,.\n\\end{equation}\nSince \n$ {\\cal D}_\\omega [ \\partial_\\varphi \\theta_0(\\varphi)]^{-T} \n= - [ \\partial_\\varphi \\theta_0(\\varphi)]^{-T} \n\\big( {\\cal D}_\\omega [\\partial_\\varphi \\theta_0(\\varphi)]^T \\big) [\\partial_\\varphi \\theta_0(\\varphi)]^{-T} $, \nthe bounds \\eqref{bella trovata 2}, \\eqref{stima A ij},\n\\eqref{ansatz 0} imply\n\\begin{equation}\\label{secondo pezzo D omega y 0 - y delta}\n\\| ({\\cal D}_\\omega [ \\partial_\\varphi \\theta_0]^{-T} ) \\rho \\|_s \n\\leq_s \\varepsilon^{6-2b} \\gamma^{-1} \n\\big( \\| Z\\|_{s + \\sigma} \\| {\\mathfrak I}_0\\|_{s_0 + \\sigma} + \\| Z \\|_{s_0 + \\sigma} \\| {\\mathfrak I}_0\\|_{s + \\sigma}\\big) \\, . \n\\end{equation}\nIn conclusion \\eqref{F diff}, \\eqref{XP delta - XP}, \\eqref{D omega y 0 - y delta},\n\\eqref{primo pezzo D omega y 0 - y delta}, \\eqref{secondo pezzo D omega y 0 - y delta} imply \n\\eqref{stima toro modificato}. The bound \\eqref{derivata i delta} follows by \n\\eqref{y 0 - y delta}, \\eqref{def Akj}, \\eqref{coefficienti pull back di Lambda}, \\eqref{ansatz 0}.\n\\end{proof}\n\nNote that there is no $ \\gamma^{- 1} $ in the right hand side of \\eqref{stima toro modificato}.\nIt turns out that an approximate inverse of $d_{i, \\zeta} {\\cal F}(i_\\delta )$\n is an approximate inverse of $d_{i, \\zeta} {\\cal F}(i_0 )$ as well. \nIn order to find an approximate inverse of the linearized operator $d_{i, \\zeta} {\\cal F}(i_\\delta )$ \nwe introduce a suitable set of symplectic coordinates nearby the isotropic torus $ i_\\d $. \nWe consider the map\n$ G_\\delta : (\\psi, \\eta, w) \\to (\\theta, y, z)$ of the phase space $\\mathbb T^\\nu \\times \\mathbb R^\\nu \\times H_S^\\bot$ defined by\n\\begin{equation}\\label{trasformazione modificata simplettica}\n\\begin{pmatrix}\n\\theta \\\\\ny \\\\\nz\n\\end{pmatrix} := G_\\delta \\begin{pmatrix}\n\\psi \\\\\n\\eta \\\\\nw\n\\end{pmatrix} := \n\\begin{pmatrix}\n\\theta_0(\\psi) \\\\\ny_\\delta (\\psi) + [\\partial_\\psi \\theta_0(\\psi)]^{-T} \\eta + \\big[ (\\partial_\\theta \\tilde{z}_0) (\\theta_0(\\psi)) \\big]^T \\partial_x^{-1} w \\\\\nz_0(\\psi) + w\n\n\\end{pmatrix} \n\\end{equation}\nwhere $\\tilde{z}_0 (\\theta) := z_0 (\\theta_0^{-1} (\\theta))$. \nIt is proved in \\cite{BB13} that $ G_\\delta $ is symplectic, using that the torus $ i_\\d $ is isotropic \n(Lemma \\ref{toro isotropico modificato}).\nIn the new coordinates, $ i_\\delta $ is the trivial embedded torus\n$ (\\psi , \\eta , w ) = (\\psi , 0, 0 ) $. \nThe transformed Hamiltonian $ K := K(\\psi, \\eta, w, \\zeta_0) $\nis (recall \\eqref{hamiltoniana modificata})\n\\begin{align} \nK := H_{\\varepsilon, \\zeta_0} \\circ G_\\d \n& = \\theta_0(\\psi) \\cdot \\zeta_0 + K_{00}(\\psi) + K_{10}(\\psi) \\cdot \\eta + (K_{0 1}(\\psi), w)_{L^2(\\mathbb T)} + \n\\frac12 K_{2 0}(\\psi)\\eta \\cdot \\eta \n\\nonumber \\\\ & \n\\quad + \\big( K_{11}(\\psi) \\eta , w \\big)_{L^2(\\mathbb T)} \n+ \\frac12 \\big(K_{02}(\\psi) w , w \\big)_{L^2(\\mathbb T)} + K_{\\geq 3}(\\psi, \\eta, w) \n\\label{KHG}\n\\end{align}\nwhere $ K_{\\geq 3} $ collects the terms at least cubic in the variables $ (\\eta, w )$.\nAt any fixed $\\psi $, \nthe Taylor coefficient $K_{00}(\\psi) \\in \\mathbb R $, \n$K_{10}(\\psi) \\in \\mathbb R^\\nu $, \n$K_{01}(\\psi) \\in H_S^\\bot$ (it is a function of $ x \\in \\mathbb T $),\n$K_{20}(\\psi) $ is a $\\nu \\times \\nu$ real matrix, \n$K_{02}(\\psi)$ is a linear self-adjoint operator of $ H_S^\\bot $ and \n$K_{11}(\\psi) : \\mathbb R^\\nu \\to H_S^\\bot$. Note that the above Taylor coefficients do not\ndepend on the parameter $ \\zeta_0 $.\n\nThe Hamilton equations associated to \\eqref{KHG} are \n\\begin{equation}\\label{sistema dopo trasformazione inverso approssimato}\n\\begin{cases}\n\\dot \\psi \\hspace{-30pt} & = K_{10}(\\psi) + K_{20}(\\psi) \\eta + \nK_{11}^T (\\psi) w + \\partial_{\\eta} K_{\\geq 3}(\\psi, \\eta, w)\n\\\\\n\\dot \\eta \\hspace{-30pt} & =- [\\partial_\\psi \\theta_0(\\psi)]^T \\zeta_0 - \n\\partial_\\psi K_{00}(\\psi) - [\\partial_{\\psi}K_{10}(\\psi)]^T \\eta - \n[\\partial_{\\psi} K_{01}(\\psi)]^T w \n\\\\\n& \\quad -\n\\partial_\\psi \\big( \\frac12 K_{2 0}(\\psi)\\eta \\cdot \\eta + ( K_{11}(\\psi) \\eta , w )_{L^2(\\mathbb T)} + \n\\frac12 ( K_{02}(\\psi) w , w )_{L^2(\\mathbb T)} + K_{\\geq 3}(\\psi, \\eta, w) \\big)\n\\\\\n\\dot w \\hspace{-30pt} & = \\partial_x \\big( K_{01}(\\psi) + \nK_{11}(\\psi) \\eta + K_{0 2}(\\psi) w + \\nabla_w K_{\\geq 3}(\\psi, \\eta, w) \\big) \n\\end{cases} \n\\end{equation}\nwhere $ [\\partial_{\\psi}K_{10}(\\psi)]^T $ is the $ \\nu \\times \\nu $ transposed matrix and \n$ [\\partial_{\\psi}K_{01}(\\psi)]^T $, $ K_{11}^T(\\psi) : {H_S^\\bot \\to \\mathbb R^\\nu} $ are defined by the \nduality relation $ ( \\partial_{\\psi} K_{01}(\\psi) [\\hat \\psi ], w)_{L^2} = \\hat \\psi \\cdot [\\partial_{\\psi}K_{01}(\\psi)]^T w $,\n$ \\forall \\hat \\psi \\in \\mathbb R^\\nu, w \\in H_S^\\bot $, \nand similarly for $ K_{11} $. \nExplicitly, for all $ w \\in H_S^\\bot $, \nand denoting $\\underline{e}_k$ the $k$-th versor of $\\mathbb R^\\nu$, \n\\begin{equation} \\label{K11 tras}\nK_{11}^T(\\psi) w = {\\mathop \\sum}_{k=1}^\\nu \\big(K_{11}^T(\\psi) w \\cdot \\underline{e}_k\\big) \\underline{e}_k =\n{\\mathop \\sum}_{k=1}^\\nu \n\\big( w, K_{11}(\\psi) \\underline{e}_k \\big)_{L^2(\\mathbb T)} \\underline{e}_k \\, \\in \\mathbb R^\\nu \\, . \n\\end{equation}\n\n\nIn the next lemma we estimate the coefficients $ K_{00} $, $ K_{10} $, $K_{01} $ \nin the Taylor expansion \\eqref{KHG}.\nNote that on an exact solution we have $ Z = 0 $ and therefore\n$ K_{00} (\\psi) = {\\rm const} $, $ K_{10} = \\omega $ and $ K_{01} = 0 $. \n\n\\begin{lemma} \\label{coefficienti nuovi}\nAssume \\eqref{ansatz 0}. Then there is $ \\sigma := \\sigma(\\tau, \\nu)$ such that \n\\[ \n\\| \\partial_\\psi K_{00} \\|_s^{{\\mathrm{Lip}(\\g)}} \n+ \\| K_{10} - \\omega \\|_s^{{\\mathrm{Lip}(\\g)}} + \\| K_{0 1} \\|_s^{{\\mathrm{Lip}(\\g)}} \n\\leq_s \\| Z \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} + \\| Z \\|_{s_0 + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\| {\\mathfrak I}_0 \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}}\\,.\n\\]\n\\end{lemma}\n\n\n\\begin{proof}\nLet $ {\\cal F}(i_\\d, \\zeta_0) := Z_\\d := (Z_{1,\\d}, Z_{2,\\d}, Z_{3,\\d}) $. \nBy a direct calculation as in \\cite{BB13} (using \\eqref{KHG}, \\eqref{operatorF})\n\\begin{align*}\n\\partial_\\psi K_{00}(\\psi) & = \n- [ \\partial_\\psi \\theta_0 (\\psi) ]^T \\big( - Z_{2, \\d} - \n[ \\partial_\\psi y_\\d] [ \\partial_\\psi \\theta_0]^{-1} Z_{1, \\d} \n+ [ (\\partial_\\theta {\\tilde z}_0)( \\theta_0 (\\psi)) ]^T \\partial_x^{-1} Z_{3,\\d} \n\\\\ \n& \\quad + [ (\\partial_\\theta {\\tilde z}_0)(\\theta_0 (\\psi)) ]^T \\partial_x^{-1} \\partial_\\psi z_0 (\\psi) [ \\partial_\\psi \\theta_0 (\\psi)]^{-1} Z_{1,\\d} \\big) \\, , \n\\\\\nK_{10}(\\psi) &\n= \\omega - [ \\partial_\\psi \\theta_0(\\psi)]^{-1} Z_{1,\\d}(\\psi) \\,, \n\\\\\nK_{01}(\\psi) \n& = - \\partial_x^{-1} Z_{3,\\d} + \\partial_x^{-1} \\partial_\\psi z_0(\\psi) [\\partial_\\psi \\theta_0(\\psi)]^{-1} Z_{1,\\d}(\\psi)\\,. \n\\end{align*}\nThen \\eqref{ansatz 0}, \\eqref{stima y - y delta}, \\eqref{stima toro modificato} (using Lemma \\ref{lemma:utile}) imply the lemma.\n\\end{proof}\n\n\\begin{remark} \\label{rem:KAM normal form} \nIf $ {\\cal F} (i_0, \\zeta_0) = 0 $ then $\\zeta_0 = 0$ by \nLemma \\ref{zeta = 0}, and Lemma \\ref{coefficienti nuovi} implies that\n\\eqref{KHG} simplifies to \n$ K = const + \\omega \\cdot \\eta + \\frac12 K_{2 0}(\\psi)\\eta \\cdot \\eta \n+ \\big( K_{11}(\\psi) \\eta , w \\big)_{L^2(\\mathbb T)} \n+ \\frac12 \\big(K_{02}(\\psi) w , w \\big)_{L^2(\\mathbb T)} + K_{\\geq 3} $. \n\\end{remark}\n\nWe now estimate $ K_{20}, K_{11}$ in \\eqref{KHG}. \nThe norm of $K_{20}$ is the sum of the norms of its matrix entries. \n\n\\begin{lemma} \\label{lemma:Kapponi vari}\nAssume \\eqref{ansatz 0}. Then \n\\begin{align}\\label{stime coefficienti K 20 11 bassa}\n\\|K_{20} + 3 \\varepsilon^{2 b} I \\|_s^{{\\mathrm{Lip}(\\g)}} \n& \\leq_s \\varepsilon^{2b+2} + \\varepsilon^{2b} \\| {\\mathfrak I}_0\\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} + \\varepsilon^{3} \\gamma^{-1} \\| {\\mathfrak I}_0\\|_{s_0 + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\| Z \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} \n\\\\ \n\\label{stime coefficienti K 11 alta} \n\\| K_{11} \\eta \\|_s^{{\\mathrm{Lip}(\\g)}} \n& \\leq_s \\varepsilon^{5} \\gamma^{-1} \\| \\eta \\|_s^{{\\mathrm{Lip}(\\g)}} \n+ \\varepsilon^{2 b - 1} ( \\| {\\mathfrak I}_0\\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} + \\gamma^{-1} \\| {\\mathfrak I}_0\\|_{s_0 + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\| Z \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} ) \n\\| \\eta \\|_{s_0}^{{\\mathrm{Lip}(\\g)}} \n\\\\\n\\label{stime coefficienti K 11 alta trasposto} \n\\| K_{11}^T w \\|_s^{{\\mathrm{Lip}(\\g)}} \n& \\leq_s \\varepsilon^{5} \\gamma^{-1} \\| w \\|_{s + 2}^{{\\mathrm{Lip}(\\g)}} \n+ \\varepsilon^{2 b - 1} ( \\| {\\mathfrak I}_0\\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} + \\gamma^{-1} \\| {\\mathfrak I}_0\\|_{s_0 + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\| Z \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} ) \n\\| w \\|_{s_0 + 2}^{{\\mathrm{Lip}(\\g)}} \\, . \n\\end{align}\nIn particular \n$ \\| K_{20} + 3 \\varepsilon^{2 b} I \\|_{s_0 }^{{\\mathrm{Lip}(\\g)}} \\leq C \\varepsilon^{6} \\gamma^{-1} $, and \n$$\n \\| K_{11} \\eta \\|_{s_0 }^{{\\mathrm{Lip}(\\g)}} \\leq C \\varepsilon^{5} \\gamma^{-1}\\| \\eta \\|_{s_0}^{{\\mathrm{Lip}(\\g)}} , \n \\quad \\| K_{11}^T w \\|_{s_0 }^{{\\mathrm{Lip}(\\g)}} \\leq C \\varepsilon^{5} \\gamma^{-1} \\| w \\|_{s_0}^{{\\mathrm{Lip}(\\g)}} \\, .\n$$ \n\\end{lemma}\n\n\\begin{proof}\nTo shorten the notation, in this proof we write $\\| \\ \\|_s$ for $\\| \\ \\|_s^{{\\mathrm{Lip}(\\g)}}$.\nWe have\n$$\nK_{2 0}(\\varphi) = [\\partial_\\varphi \\theta_0(\\varphi)]^{-1} \\partial_{yy} H_\\varepsilon(i_\\delta(\\varphi)) \n[\\partial_\\varphi \\theta_0(\\varphi)]^{-T} \n= [\\partial_\\varphi \\theta_0(\\varphi)]^{-1} \\partial_{yy} P(i_\\delta(\\varphi)) \n[\\partial_\\varphi \\theta_0(\\varphi)]^{-T}.\n$$\nThen \\eqref{D yy P}, \\eqref{ansatz 0}, \\eqref{stima y - y delta} \nimply \\eqref{stime coefficienti K 20 11 bassa}.\nNow (see also \\cite{BB13})\n\\begin{align*}\nK_{11}(\\varphi) & = \n\\partial_{y} \\nabla_z H_\\varepsilon (i_\\delta(\\varphi)) [\\partial_\\varphi \\theta_0 (\\varphi)]^{-T} \n- \\partial_x^{-1} (\\partial_\\theta {\\tilde z}_0) (\\theta_0(\\varphi)) (\\partial_{yy} H_\\varepsilon) (i_\\delta(\\varphi)) [\\partial_\\varphi \\theta_0 (\\varphi)]^{-T} \\nonumber\\\\\n& \\stackrel{\\eqref{Hamiltoniana Heps KdV}} = \n \\partial_{y} \\nabla_z P(i_\\delta(\\varphi)) [\\partial_\\varphi \\theta_0 (\\varphi)]^{-T} \n- \\partial_x^{-1} (\\partial_\\theta {\\tilde z}_0) (\\theta_0(\\varphi)) (\\partial_{yy} P) (i_\\delta(\\varphi)) [\\partial_\\varphi \\theta_0 (\\varphi)]^{-T}\\,,\n\\end{align*}\ntherefore, using \\eqref{stime linearizzato campo hamiltoniano}, \\eqref{D yy P}, \\eqref{ansatz 0},\nwe deduce \\eqref{stime coefficienti K 11 alta}. \nThe bound \\eqref{stime coefficienti K 11 alta trasposto} for $K_{11}^T$ follows by \\eqref{K11 tras}.\n\\end{proof}\n\nUnder the linear change of variables \n\\begin{equation}\\label{DGdelta}\nD G_\\delta(\\varphi, 0, 0) \n\\begin{pmatrix}\n\\widehat \\psi \\, \\\\\n\\widehat \\eta \\\\\n\\widehat w\n\\end{pmatrix} \n:= \n\\begin{pmatrix}\n\\partial_\\psi \\theta_0(\\varphi) & 0 & 0 \\\\\n\\partial_\\psi y_\\delta(\\varphi) & [\\partial_\\psi \\theta_0(\\varphi)]^{-T} & \n- [(\\partial_\\theta \\tilde{z}_0)(\\theta_0(\\varphi))]^T \\partial_x^{-1} \\\\\n\\partial_\\psi z_0(\\varphi) & 0 & I\n\\end{pmatrix}\n\\begin{pmatrix}\n\\widehat \\psi \\, \\\\\n\\widehat \\eta \\\\\n\\widehat w\n\\end{pmatrix} \n\\end{equation}\nthe linearized operator $d_{i, \\zeta}{\\cal F}(i_\\delta )$ \ntransforms (approximately, see \\eqref{verona 2}) into the operator obtained linearizing \n\\eqref{sistema dopo trasformazione inverso approssimato} at $(\\psi, \\eta , w, \\zeta ) = (\\varphi, 0, 0, \\zeta_0 )$\n(with $ \\partial_t \\rightsquigarrow {\\cal D}_\\omega $), namely \n\\begin{equation}\\label{lin idelta}\n\\hspace{-5pt}\n\\begin{pmatrix}\n{\\cal D}_\\omega \\widehat \\psi - \\partial_\\psi K_{10}(\\varphi)[\\widehat \\psi \\, ] - \nK_{2 0}(\\varphi)\\widehat \\eta - K_{11}^T (\\varphi) \\widehat w \\\\\n {\\cal D}_\\omega \\widehat \\eta + [\\partial_\\psi \\theta_0(\\varphi)]^T \\widehat \\zeta + \n\\partial_\\psi [\\partial_\\psi \\theta_0(\\varphi)]^T [ \\widehat \\psi, \\zeta_0] + \\partial_{\\psi\\psi} K_{00}(\\varphi)[\\widehat \\psi] + \n[\\partial_\\psi K_{10}(\\varphi)]^T \\widehat \\eta + \n[\\partial_\\psi K_{01}(\\varphi)]^T \\widehat w \\\\ \n{\\cal D}_\\omega \\widehat w - \\partial_x \n\\{ \\partial_\\psi K_{01}(\\varphi)[\\widehat \\psi] + K_{11}(\\varphi) \\widehat \\eta + K_{02}(\\varphi) \\widehat w \\}\n\\end{pmatrix} \\! . \\hspace{-5pt}\n\\end{equation}\nWe now estimate the induced composition operator. \n\\begin{lemma} \\label{lemma:DG}\nAssume \\eqref{ansatz 0} and let \n$ \\widehat \\imath := (\\widehat \\psi, \\widehat \\eta, \\widehat w)$. \nThen \n\\begin{gather} \\label{DG delta}\n\\|DG_\\delta(\\varphi,0,0) [\\widehat \\imath] \\|_s + \\|DG_\\delta(\\varphi,0,0)^{-1} [\\widehat \\imath] \\|_s \n\\leq_s \\| \\widehat \\imath \\|_{s} + ( \\| {\\mathfrak I}_0 \\|_{s + \\sigma} + \n\\gamma^{-1} \\| {\\mathfrak I}_0 \\|_{s_0 + \\sigma} \\|Z \\|_{s + \\sigma} ) \\| \\widehat \\imath \\|_{s_0}\\,,\n\\\\ \n\\| D^2 G_\\delta(\\varphi,0,0)[\\widehat \\imath_1, \\widehat \\imath_2] \\|_s \n\\leq_s \\| \\widehat \\imath_1\\|_s \\| \\widehat \\imath_2 \\|_{s_0} \n+ \\| \\widehat \\imath_1\\|_{s_0} \\| \\widehat \\imath_2 \\|_{s} \n+ ( \\| {\\mathfrak I}_0 \\|_{s + \\sigma} + \\gamma^{-1} \\| {\\mathfrak I}_0 \\|_{s_0 + \\sigma} \\| Z\\|_{s + \\sigma} ) \\|\\widehat \\imath_1 \\|_{s_0} \\| \\widehat \\imath_2\\|_{s_0}\n\\notag \n\\end{gather}\nfor some $\\sigma := \\sigma(\\nu,\\t)$. \nMoreover the same estimates hold if we replace the norm $\\| \\ \\|_s$ with $\\| \\ \\|_s^{{\\mathrm{Lip}(\\g)}}$.\n\\end{lemma}\n\n\n\\begin{proof}\nThe estimate \\eqref{DG delta} for $D G_\\delta(\\varphi,0,0)$ follows by \\eqref{DGdelta} and \\eqref{stima y - y delta}.\nBy \\eqref{ansatz 0},\n$ \\| (DG_\\delta(\\varphi,0,0) - I) \\widehat \\imath \\|_{s_0} \n\\leq $ $ C \\varepsilon^{6 - 2b} \\gamma^{-1} \\| \\widehat \\imath \\|_{s_0} \n\\leq \\| \\widehat \\imath \\|_{s_0} \/ 2 $.\nTherefore $DG_\\delta(\\varphi,0,0)$ is invertible and, by Neumann series, the inverse \nsatisfies \\eqref{DG delta}. \nThe bound for $D^2 G_\\delta$ follows by differentiating $DG_\\d$. \n\\end{proof}\n\n\nIn order to construct an approximate inverse of \\eqref{lin idelta} it is sufficient to solve the equation\n\\begin{equation}\\label{operatore inverso approssimato} \n{\\mathbb D} [\\widehat \\psi, \\widehat \\eta, \\widehat w, \\widehat \\zeta ] := \n \\begin{pmatrix}\n{\\cal D}_\\omega \\widehat \\psi - \nK_{20}(\\varphi) \\widehat \\eta - K_{11}^T(\\varphi) \\widehat w\\\\\n{\\cal D}_\\omega \\widehat \\eta + [\\partial_\\psi \\theta_0(\\varphi)]^T \\widehat \\zeta \\\\\n{\\cal D}_\\omega \\widehat w - \\partial_x K_{11}(\\varphi)\\widehat \\eta - \\partial_x K_{0 2}(\\varphi) \\widehat w \n\\end{pmatrix} =\n\\begin{pmatrix}\ng_1 \\\\ g_2 \\\\ g_3 \n\\end{pmatrix}\n\\end{equation}\nwhich is obtained by neglecting in \\eqref{lin idelta} \nthe terms $ \\partial_\\psi K_{10} $, $ \\partial_{\\psi \\psi} K_{00} $, $ \\partial_\\psi K_{00} $, $ \\partial_\\psi K_{01} $ and\n$ \\partial_\\psi [\\partial_\\psi \\theta_0(\\varphi)]^T [ \\cdot , \\zeta_0] $\n(which are naught at a solution by Lemmata \\ref{coefficienti nuovi} and \\ref{zeta = 0}). \n\nFirst we solve the second equation in \\eqref{operatore inverso approssimato}, namely \n$ {\\cal D}_\\omega \\widehat \\eta = g_2 - [\\partial_\\psi \\theta_0(\\varphi)]^T \\widehat \\zeta $. \nWe choose $ \\widehat \\zeta $ so that the $\\varphi$-average of the right hand side \nis zero, namely \n\\begin{equation}\\label{fisso valore di widehat zeta}\n\\widehat \\zeta = \\langle g_2 \\rangle \n\\end{equation}\n(we denote $ \\langle g \\rangle := (2 \\pi)^{- \\nu} \\int_{\\mathbb T^\\nu} g (\\varphi) d \\varphi $). \nNote that the $\\varphi$-averaged matrix $ \\langle [\\partial_\\psi \\theta_0 ]^T \\rangle = \n\\langle I + [\\partial_\\psi \\Theta_0]^T \\rangle = I $ because $\\theta_0(\\varphi) = \\varphi + \\Theta_0(\\varphi)$ and \n$\\Theta_0(\\varphi)$ is a periodic function.\nTherefore \n\\begin{equation}\\label{soleta}\n\\widehat \\eta := {\\cal D}_\\omega^{-1} \\big(\ng_2 - [\\partial_\\psi \\theta_0(\\varphi) ]^T \\langle g_2 \\rangle \\big) + \n\\langle \\widehat \\eta \\rangle \\, , \\quad \\langle \\widehat \\eta \\rangle \\in \\mathbb R^\\nu \\, , \n\\end{equation}\nwhere the average $\\langle \\widehat \\eta \\rangle$ will be fixed below. \nThen we consider the third equation\n\\begin{equation}\\label{cal L omega}\n{\\cal L}_\\omega \\widehat w = g_3 + \\partial_x K_{11}(\\varphi) \\widehat \\eta\\,, \\\n\\quad {\\cal L}_\\omega := \\omega \\cdot \\partial_\\varphi - \\partial_x K_{0 2}(\\varphi) \\, .\n\\end{equation}\n\n\\begin{itemize}\n\\item {\\sc Inversion assumption.} \n{\\it There exists a set $ \\Omega_\\infty \\subset \\Omega_o$ such that \nfor all $ \\omega \\in \\Omega_\\infty $, for every function\n$ g \\in H^{s+\\mu}_{S^\\bot} (\\mathbb T^{\\nu+1}) $ there \nexists a solution $ h := {\\cal L}_\\omega^{- 1} g \\in H^{s}_{S^\\bot} (\\mathbb T^{\\nu+1}) $ \nof the linear equation $ {\\cal L}_\\omega h = g $\nwhich satisfies}\n\\begin{equation}\\label{tame inverse}\n\\| {\\cal L}_\\omega^{- 1} g \\|_s^{{\\mathrm{Lip}(\\g)}} \\leq C(s) \\gamma^{-1} \n\\big( \\| g \\|_{s + \\mu}^{{\\mathrm{Lip}(\\g)}} + \\varepsilon \\gamma^{-1} \n\\big\\{ \\| {\\mathfrak I}_0 \\|_{s + \\mu}^{{\\mathrm{Lip}(\\g)}} + \\gamma^{-1} \\| {\\mathfrak I}_0 \\|_{s_0 + \\mu}^{{\\mathrm{Lip}(\\g)}} \\| Z \\|_{s + \\mu}^{{\\mathrm{Lip}(\\g)}} \\big\\} \\|g \\|_{s_0}^{{\\mathrm{Lip}(\\g)}} \\big) \n\\end{equation}\n\\emph{for some $ \\mu := \\mu (\\tau, \\nu) > 0 $}. \n\\end{itemize}\n\n\n\\begin{remark}\nThe term $ \\varepsilon \\gamma^{-1} \n\\{ \\| {\\mathfrak I}_0 \\|_{s + \\mu}^{{\\mathrm{Lip}(\\g)}} + \\gamma^{-1} \\| {\\mathfrak I}_0 \\|_{s_0 + \\mu}^{{\\mathrm{Lip}(\\g)}} \\| Z \\|_{s + \\mu}^{{\\mathrm{Lip}(\\g)}} \\} $ arises because\nthe remainder $ R_6 $ in section \\ref{step5} contains the term \n$ \\varepsilon ( \\| \\Theta_0 \\|_{s + \\mu}^{{\\mathrm{Lip}(\\g)}} + \\| y_\\d \\|_{s + \\mu}^{{\\mathrm{Lip}(\\g)}}) $ \n$\\leq \\varepsilon \\| {\\mathfrak I}_\\d \\|_{s + \\mu}^{{\\mathrm{Lip}(\\g)}} $, see Lemma \\ref{lemma L6}. \n\\end{remark}\n\nBy the above assumption there exists a solution \n\\begin{equation}\\label{normalw}\n\\widehat w := {\\cal L}_\\omega^{-1} [ g_3 + \\partial_x K_{11}(\\varphi) \\widehat \\eta \\, ] \n\\end{equation}\nof \\eqref{cal L omega}. \nFinally, we solve the first equation in \\eqref{operatore inverso approssimato}, \nwhich, substituting \\eqref{soleta}, \\eqref{normalw}, becomes\n\\begin{equation}\\label{equazione psi hat}\n{\\cal D}_\\omega \\widehat \\psi = \ng_1 + M_1(\\varphi) \\langle \\widehat \\eta \\rangle + M_2(\\varphi) g_2 + M_3(\\varphi) g_3 - \nM_2(\\varphi)[\\partial_\\psi \\theta_0]^T \\langle g_2 \\rangle \\,,\n\\end{equation}\nwhere\n\\begin{equation} \\label{cal M2}\nM_1(\\varphi) := K_{2 0}(\\varphi) + K_{11}^T(\\varphi) {\\cal L}_\\omega^{-1} \\partial_x K_{11}(\\varphi)\\,, \\quad\nM_2(\\varphi) := M_1 (\\varphi) {\\cal D}_\\omega^{-1} \\, , \\quad \nM_3(\\varphi) := K_{11}^T (\\varphi) {\\cal L}_\\omega^{-1} \\, . \n\\end{equation}\nIn order to solve the equation \\eqref{equazione psi hat} we have \nto choose $\\langle \\widehat \\eta \\rangle$ such that the right hand side in \\eqref{equazione psi hat} has zero average. \nBy Lemma \\ref{lemma:Kapponi vari} and \\eqref{ansatz 0}, the $\\varphi$-averaged matrix \n$ \\langle M_1 \\rangle =- 3 \\varepsilon^{2 b} I + O( \\varepsilon^{10} \\gamma^{-3}) $. \nTherefore, for $ \\varepsilon $ small, $\\langle M_1 \\rangle$ is invertible and $\\langle M_1 \\rangle^{-1} = O(\\varepsilon^{-2 b}) = O(\\gamma^{- 1})$ \n(recall \\eqref{link gamma b}). Thus we define \n\\begin{equation}\\label{sol alpha}\n\\langle \\widehat \\eta \\rangle := - \\langle M_1 \\rangle^{-1} \n[ \\langle g_1 \\rangle + \\langle M_2 g_2 \\rangle + \\langle M_3 g_3 \\rangle - \n\\langle M_2 [\\partial_\\psi \\theta_0]^T \\rangle \\langle g_2 \\rangle ].\n\\end{equation}\nWith this choice of $\\langle \\widehat \\eta \\rangle$ the equation \\eqref{equazione psi hat} has the solution\n\\begin{equation}\\label{sol psi}\n\\widehat \\psi :=\n{\\cal D}_\\omega^{-1} [ g_1 + M_1(\\varphi) \\langle \\widehat \\eta \\rangle + M_2(\\varphi) g_2 + M_3(\\varphi) g_3 -\nM_2(\\varphi)[\\partial_\\psi \\theta_0]^T \\langle g_2 \\rangle ].\n\\end{equation}\nIn conclusion, we have constructed \na solution $(\\widehat \\psi, \\widehat \\eta, \\widehat w, \\widehat \\zeta)$ of the linear system \\eqref{operatore inverso approssimato}. \n\n\\begin{proposition}\\label{prop: ai}\nAssume \\eqref{ansatz 0} and \\eqref{tame inverse}. \nThen, $\\forall \\omega \\in \\Omega_\\infty $, $ \\forall g := (g_1, g_2, g_3) $,\n the system \\eqref{operatore inverso approssimato} has a solution \n$ {\\mathbb D}^{-1} g := (\\widehat \\psi, \\widehat \\eta, \\widehat w, \\widehat \\zeta ) $\nwhere $(\\widehat \\psi, \\widehat \\eta, \\widehat w, \\widehat \\zeta)$ are defined in \n\\eqref{sol psi}, \\eqref{soleta}, \\eqref{sol alpha}, \\eqref{normalw}, \\eqref{fisso valore di widehat zeta} satisfying\n\\begin{equation} \\label{stima T 0 b}\n\\| {\\mathbb D}^{-1} g \\|_s^{{\\rm Lip}(\\gamma)} \n\\leq_s \\gamma^{-1} \\big( \\| g \\|_{s + \\mu}^{{\\rm Lip}(\\gamma)} \n+ \\varepsilon \\gamma^{-1} \\big\\{ \\| {\\mathfrak I}_0 \\|_{s + \\mu}^{{\\rm Lip}(\\gamma)} \n+ \\gamma^{-1} \\| {\\mathfrak I}_0 \\|_{s_0 + \\mu}^{{\\rm Lip}(\\gamma)} \\|{\\cal F}(i_0, \\zeta_0) \\|^{{\\rm Lip}(\\gamma)}_{s + \\mu} \\big\\} \\| g \\|_{s_0 + \\mu}^{{\\rm Lip}(\\gamma)} \\big).\n\\end{equation}\n\\end{proposition}\n\n\\begin{proof}\n\nRecalling \\eqref{cal M2}, by Lemma \\ref{lemma:Kapponi vari}, \\eqref{tame inverse}, \\eqref{ansatz 0} we get $ \\| M_2 h \\|_{s_0} + \\| M_3 h \\|_{s_0} \\leq C \\| h \\|_{s_0 + \\sigma} $. \nThen, by \\eqref{sol alpha} and $\\langle M_1 \\rangle^{-1} = O(\\varepsilon^{-2 b}) = O(\\gamma^{-1}) $, \nwe deduce \n$ |\\langle \\widehat \\eta\\rangle|^{{\\mathrm{Lip}(\\g)}} \\leq C\\gamma^{-1} \\| g \\|_{s_0+ \\sigma}^{{\\mathrm{Lip}(\\g)}} $\nand \\eqref{soleta}, \\eqref{Dom inverso} imply \n$ \\| \\widehat \\eta \\|_s^{{\\mathrm{Lip}(\\g)}} \\leq_s \\gamma^{-1} \\big( \\| g \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} + \\| {\\mathfrak{I}}_0 \\|_{s + \\sigma } \\| g \\|_{s_0}^{\\mathrm{Lip}(\\g)} \\big)$.\nThe bound \\eqref{stima T 0 b} is sharp for $ \\widehat w $ because $ {\\cal L}_\\omega^{-1} g_3 $ in \\eqref{normalw}\nis estimated using \\eqref{tame inverse}. Finally $ \\widehat \\psi $ \nsatisfies \\eqref{stima T 0 b} using\n\\eqref{sol psi}, \\eqref{cal M2}, \\eqref{tame inverse}, \\eqref{Dom inverso} and Lemma \\ref{lemma:Kapponi vari}.\n\\end{proof}\nFinally we prove that the operator \n\\begin{equation}\\label{definizione T} \n{\\bf T}_0 := (D { \\widetilde G}_\\delta)(\\varphi,0,0) \\circ {\\mathbb D}^{-1} \\circ (D G_\\delta) (\\varphi,0,0)^{-1}\n\\end{equation}\nis an approximate right inverse for $d_{i,\\zeta} {\\cal F}(i_0 )$ where\n$ \\widetilde{G}_\\delta (\\psi, \\eta, w, \\zeta) := $ $ \\big( G_\\delta (\\psi, \\eta, w), \\zeta \\big) $ \n is the identity on the $ \\zeta $-component. \nWe denote the norm $ \\| (\\psi, \\eta, w, \\zeta) \\|_s^{\\mathrm{Lip}(\\g)} := $ $ \\max \\{ \\| (\\psi, \\eta, w) \\|_s^{\\mathrm{Lip}(\\g)}, $ $ | \\zeta |^{\\mathrm{Lip}(\\g)} \\} $.\n\n\n\\begin{theorem} {\\bf (Approximate inverse)} \\label{thm:stima inverso approssimato}\nAssume \\eqref{ansatz 0} and the inversion assumption \\eqref{tame inverse}. \nThen there exists $ \\mu := \\mu (\\tau, \\nu) > 0 $ such that, for all $ \\omega \\in \\Omega_\\infty $, \nfor all $ g := (g_1, g_2, g_3) $, \nthe operator $ {\\bf T}_0 $ defined in \\eqref{definizione T} satisfies \n\\begin{equation}\\label{stima inverso approssimato 1}\n\\| {\\bf T}_0 g \\|_{s}^{{\\rm Lip}(\\gamma)} \n\\leq_s \\gamma^{-1} \\big(\\| g \\|_{s + \\mu}^{{\\rm Lip}(\\gamma)} \n+ \\varepsilon \\gamma^{-1} \\big\\{ \\| {\\mathfrak I}_0 \\|_{s + \\mu}^{{\\rm Lip}(\\gamma)} \n+\\gamma^{-1} \\| {\\mathfrak{I}}_0 \\|_{s_0 + \\mu}^{\\mathrm{Lip}(\\g)} \n\\|{\\cal F}(i_0, \\zeta_0) \\|_{s + \\mu}^{{\\rm Lip}(\\gamma)} \\big\\} \n\\| g \\|_{s_0 + \\mu}^{{\\rm Lip}(\\gamma)} \\big). \n\\end{equation}\nIt is an approximate inverse of $d_{i, \\zeta} {\\cal F}(i_0 )$, namely \n\\begin{align}\n& \\| ( d_{i, \\zeta} {\\cal F}(i_0) \\circ {\\bf T}_0 - I ) g \\|_s^{{\\rm Lip}(\\gamma)} \n\\label{stima inverso approssimato 2} \n\\\\ \n& \\leq_s \\gamma^{-1} \\Big( \\| {\\cal F}(i_0, \\zeta_0) \\|_{s_0 + \\mu}^{\\mathrm{Lip}(\\g)} \\| g \\|_{s + \\mu}^{\\mathrm{Lip}(\\g)} \n+ \\big\\{ \\| {\\cal F}(i_0, \\zeta_0) \\|_{s + \\mu}^{\\mathrm{Lip}(\\g)} \n+ \\varepsilon \\gamma^{-1} \\| {\\cal F}(i_0, \\zeta_0) \\|_{s_0 + \\mu}^{\\mathrm{Lip}(\\g)} \\| {\\mathfrak I}_0 \\|_{s + \\mu}^{\\mathrm{Lip}(\\g)} \\big\\} \\| g \\|_{s_0 + \\mu}^{\\mathrm{Lip}(\\g)} \\Big).\n\\nonumber\n\\end{align}\n\\end{theorem}\n\n\n\\begin{proof}\nWe denote $\\| \\ \\|_s$ instead of $\\| \\ \\|_s^{{\\mathrm{Lip}(\\g)}}$. \nThe bound \\eqref{stima inverso approssimato 1} \nfollows from \\eqref{definizione T}, \\eqref{stima T 0 b}, \\eqref{DG delta}.\nBy \\eqref{operatorF}, since $ X_\\mathcal{N} $ does not depend on $ y $, \nand $ i_\\d $ differs from $ i_0 $ only for the $ y$ component, \nwe have \n\\begin{align} \\label{verona 0}\nd_{i, \\zeta} {\\cal F}(i_0 )[\\, \\widehat \\imath, \\widehat \\zeta \\, ] - d_{i, \\zeta} {\\cal F}(i_\\delta ) [\\, \\widehat \\imath, \\widehat \\zeta \\, ]\n& = d_i X_P (i_\\delta) [\\, \\widehat \\imath \\, ] - d_i X_P (i_0) [\\, \\widehat \\imath \\, ]\n\\\\ & \n= \\int_0^1 \\partial_y d_i X_P (\\theta_0, y_0 + s (y_\\delta - y_0), z_0) [y_\\d - y_0, \\widehat \\imath \\, ] ds \n=: {\\cal E}_0 [\\, \\widehat \\imath, \\widehat \\zeta \\, ] \\,. \\nonumber\n\\end{align}\nBy \\eqref{D yii}, \\eqref{stima y - y delta}, \\eqref{ansatz 0}, we estimate\n\\begin{equation}\\label{stima parte trascurata 1}\n\\| {\\cal E}_0 [\\, \\widehat \\imath, \\widehat \\zeta \\, ] \\|_s \\leq_s \n\\| Z \\|_{s_0 + \\sigma} \\| \\widehat \\imath \\|_{s + \\sigma} + \n\\| Z \\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma} + \\varepsilon^{2b-1}\\gamma^{-1}\n\\| Z \\|_{s_0 + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma} \\| {\\mathfrak{I}}_0 \\|_{s+\\sigma} \n\\end{equation}\nwhere $Z := \\mathcal{F}(i_0, \\zeta_0)$ (recall \\eqref{def Zetone}). \n Note that $\\mathcal{E}_0[\\widehat \\imath, \\widehat \\zeta]$ is, in fact, independent of $\\widehat \\zeta$.\nDenote the set of variables $ (\\psi, \\eta, w) =: {\\mathtt u} $. \nUnder the transformation $G_\\delta $, the nonlinear operator ${\\cal F}$ in \\eqref{operatorF} transforms into \n\\begin{equation} \\label{trasfo imp}\n{\\cal F}(G_\\delta( {\\mathtt u} (\\varphi) ), \\zeta ) \n= D G_\\delta( {\\mathtt u} (\\varphi) ) \\big( {\\cal D}_\\omega {\\mathtt u} (\\varphi) - X_K ( {\\mathtt u} (\\varphi), \\zeta) \\big) \\, , \\quad \nK = H_{\\varepsilon, \\zeta} \\circ G_\\delta \\, ,\n\\end{equation}\nsee \\eqref{sistema dopo trasformazione inverso approssimato}. \nDifferentiating \\eqref{trasfo imp} at the trivial torus \n$ {\\mathtt u}_\\delta (\\varphi) = G_\\delta^{-1}(i_\\delta) (\\varphi) = (\\varphi, 0 , 0 ) $, \nat $ \\zeta = \\zeta_0 $, \nin the directions $(\\widehat {\\mathtt u}, \\widehat \\zeta)\n= (D G_\\d ({\\mathtt u}_\\d)^{-1} [\\, \\widehat \\imath \\, ], \\widehat \\zeta) \n= D {\\widetilde G}_\\d ({\\mathtt u}_\\d)^{-1} [\\, \\widehat \\imath , \\widehat \\zeta \\, ] $, \nwe get\n\\begin{align} \\label{verona 2}\nd_{i , \\zeta} {\\cal F}(i_\\delta ) [\\, \\widehat \\imath, \\widehat \\zeta \\, ]\n= & D G_\\delta( {\\mathtt u}_\\delta) \n\\big( {\\cal D}_\\omega \\widehat {\\mathtt u} \n- d_{\\mathtt u, \\zeta} X_K( {\\mathtt u}_\\delta, \\zeta_0) [\\widehat {\\mathtt u}, \\widehat \\zeta \\, ] \n\\big) \n+ {\\cal E}_1 [ \\, \\widehat \\imath , \\widehat \\zeta \\, ]\\,,\n\\\\\n\\label{E1}\n{\\cal E}_1 [\\, \\widehat \\imath , \\widehat \\zeta \\, ] \n:= & \nD^2 G_\\delta( {\\mathtt u}_\\delta) \\big[ D G_\\delta( {\\mathtt u}_\\delta)^{-1} {\\cal F}(i_\\delta, \\zeta_0), \\, D G_\\d({\\mathtt u}_\\d)^{-1} \n[ \\, \\widehat \\imath \\, ] \\big] \\,, \n\\end{align}\nwhere $ d_{\\mathtt u, \\zeta} X_K( {\\mathtt u}_\\delta, \\zeta_0) $ is expanded in \\eqref{lin idelta}.\nIn fact, ${\\cal E}_1$ is independent of $\\widehat \\zeta$. \nWe split \n\\[ \n{\\cal D}_\\omega \\widehat {\\mathtt u} \n- d_{\\mathtt u, \\zeta} X_K( {\\mathtt u}_\\delta, \\zeta_0) [\\widehat {\\mathtt u}, \\widehat \\zeta] \n= \\mathbb{D} [\\widehat {\\mathtt u}, \\widehat \\zeta \\, ] + R_Z [ \\widehat {\\mathtt u}, \\widehat \\zeta \\, ], \n\\]\nwhere $ {\\mathbb D} [\\widehat {\\mathtt u}, \\widehat \\zeta] $ is defined in \\eqref{operatore inverso approssimato} and \n\\begin{equation}\\label{R0}\nR_Z [ \\widehat \\psi, \\widehat \\eta, \\widehat w, \\widehat \\zeta]\n:= \\begin{pmatrix}\n - \\partial_\\psi K_{10}(\\varphi) [\\widehat \\psi ] \\\\\n\\partial_\\psi [\\partial_\\psi \\theta_0(\\varphi)]^T [ \\widehat \\psi, \\zeta_0] + \\partial_{\\psi \\psi} K_{00} (\\varphi) [ \\widehat \\psi ] + \n [\\partial_\\psi K_{10}(\\varphi)]^T \\widehat \\eta + \n [\\partial_\\psi K_{01}(\\varphi)]^T \\widehat w \\\\\n - \\partial_x \\{ \\partial_{\\psi} K_{01}(\\varphi)[ \\widehat \\psi ] \\}\n \\end{pmatrix}\n\\end{equation}\n($R_Z$ is independent of $\\widehat \\zeta$). \nBy \\eqref{verona 0} and \\eqref{verona 2}, \n\\begin{equation} \\label{E2}\nd_{i, \\zeta} {\\cal F}(i_0 ) \n= D G_\\delta({\\mathtt u}_\\delta) \\circ {\\mathbb D} \\circ D {\\widetilde G}_\\delta ({\\mathtt u}_\\delta)^{-1} \n+ {\\cal E}_0 + {\\cal E}_1 + \\mathcal{E}_2 \\,,\n\\quad\n\\mathcal{E}_2 := D G_\\delta( {\\mathtt u}_\\delta) \\circ R_Z \\circ D {\\widetilde G}_\\delta ({\\mathtt u}_\\delta)^{-1} \\, .\n\\end{equation}\nBy Lemmata \\ref{coefficienti nuovi}, \\ref{lemma:DG}, \\ref{zeta = 0}, and \\eqref{stima toro modificato}, \\eqref{ansatz 0},\nthe terms $\\mathcal{E}_1, \\mathcal{E}_2 $ (see \\eqref{E1}, \\eqref{E2}, \\eqref{R0}) satisfy \nthe same bound \\eqref{stima parte trascurata 1} as $\\mathcal{E}_0$\n(in fact even better).\nThus the sum $\\mathcal{E} := \\mathcal{E}_0 + \\mathcal{E}_1 + \\mathcal{E}_2$ satisfies \\eqref{stima parte trascurata 1}.\nApplying $ {\\bf T}_0 $ defined in \\eqref{definizione T} to the right in \\eqref{E2}, \nsince $ {\\mathbb D} \\circ {\\mathbb D}^{-1} = I $ (see Proposition \\ref{prop: ai}), \nwe get $d_{i, \\zeta} {\\cal F}(i_0 ) \\circ {\\bf T}_0 - I \n= \\mathcal{E} \\circ {\\bf T}_0$. \nThen \\eqref{stima inverso approssimato 2} follows from \n\\eqref{stima inverso approssimato 1} and the bound \\eqref{stima parte trascurata 1} for $\\mathcal{E}$. \n\\end{proof}\n\n\n\n\\section{The linearized operator in the normal directions}\\label{linearizzato siti normali}\n\nThe goal of this section is to write an explicit expression of the linearized operator \n$\\mathcal{L}_\\omega$ defined in \\eqref{cal L omega}, see Proposition \\ref{prop:lin}.\nTo this aim, we compute $ \\frac12 ( K_{02}(\\psi) w, w )_{L^2(\\mathbb T)} $, $ w \\in H_S^\\bot$, \nwhich collects all the components of $(H_\\varepsilon \\circ G_\\d)(\\psi, 0, w)$ that are quadratic in $w$, see \\eqref{KHG}.\n\n\\smallskip\n\nWe first prove some preliminary lemmata. \n\n\\begin{lemma}\\label{lemma astratto potente}\nLet $ H $ be a Hamiltonian of class $C^2 ( H^1_0(\\mathbb T_x), \\mathbb R )$ and consider a map\n$ \\Phi(u) := u + \\Psi(u) $ satisfying $\\Psi (u) = \\Pi_E \\Psi(\\Pi_E u)$, for all $ u $, \nwhere $E$ is a finite dimensional subspace as in \\eqref{def E finito}. Then \n\\begin{equation}\\label{lint2}\n\\partial_u \\big[\\nabla ( H \\circ \\Phi)\\big] (u) [h] = (\\partial_u \\nabla H )(\\Phi(u)) [h] + {\\cal R}(u)[h]\\,,\n\\end{equation}\nwhere $ {\\cal R}(u) $ \nhas the ``finite dimensional\" form \n\\begin{equation}\\label{forma buona resto}\n{\\cal R}(u)[h] = {\\mathop\\sum}_{|j| \\leq C} \\big( h , g_j(u) \\big)_{L^2(\\mathbb T)} \\chi_j(u) \n\\end{equation}\nwith $ \\chi_j (u) = e^{{\\mathrm i} j x} $ or $ g_j(u) = e^{{\\mathrm i} j x} $. The remainder \n$ {\\cal R} (u) = {\\cal R}_0 (u) + {\\cal R}_1 (u) + {\\cal R}_2 (u) $ with\n\\begin{align}\\label{resti012} \n{\\cal R}_0 (u) & := (\\partial_u \\nabla H)(\\Phi(u)) \\partial_u \\Psi (u), \\qquad \n{\\cal R}_1 (u) := [\\partial_{u }\\{ \\Psi'(u)^T\\}] [ \\cdot , \\nabla H(\\Phi(u)) ], \\nonumber \\\\\n\\, {\\cal R}_2 (u) & := [\\partial_u \\Psi (u)]^T (\\partial_u \\nabla H)(\\Phi(u)) \\partial_u \\Phi(u). \n\\end{align}\n\\end{lemma}\n\n\\begin{proof}\nBy a direct calculation, \n\\begin{equation}\\label{nabla composto} \n\\nabla (H \\circ \\Phi)(u) = [\\Phi'(u)]^T \\nabla H(\\Phi(u)) = \\nabla H(\\Phi(u)) + [\\Psi'(u)]^T \\nabla H(\\Phi(u))\n\\end{equation}\nwhere $ \\Phi' (u) := ( \\partial_u \\Phi) (u) $ and $ [ \\ ]^T $ denotes the transpose with respect to the $ L^2 $ scalar product.\nDifferentiating \\eqref{nabla composto}, we get \\eqref{lint2} and \\eqref{resti012}.\n\nLet us show that each $ {\\cal R}_m $ has the form \\eqref{forma buona resto}. \nWe have \n\\begin{equation}\\label{Psiuh}\n\\Psi'(u) = \\Pi_E \\Psi'(\\Pi_E u) \\Pi_E \\, \\, , \\quad [\\Psi'(u)]^T = \\Pi_E [\\Psi'( \\Pi_E u)]^T \\Pi_E \\, . \n\\end{equation}\nHence, setting $ A := (\\partial_u \\nabla H)(\\Phi(u)) \\Pi_E \\Psi' ( \\Pi_E u) $, we get \n\\[\n{\\cal R}_0(u)[h] \n= A [ \\Pi_E h ] \n= {\\mathop\\sum}_{|j| \\leq C } h_j A ( e^{{\\mathrm i} j x} ) \n= {\\mathop\\sum}_{|j| \\leq C} (h, g_j )_{L^2(\\mathbb T)} \\chi_j \n\\]\nwith $ g_j := e^{{\\mathrm i} j x} $, $ \\chi_j := A( e^{{\\mathrm i} jx } )$.\nSimilarly, using \\eqref{Psiuh}, and setting \n$ A := [ \\Psi' (\\Pi_E u) ]^T \\Pi_E (\\partial_u \\nabla H)(\\Phi(u)) \\Phi'(u) $, \nwe get\n$$\n{\\cal R}_2 (u)[h] = \\Pi_E [ A h ] = {\\mathop\\sum}_{|j | \\leq C } (A h , e^{{\\mathrm i} j x} )_{L^2(\\mathbb T)} e^{{\\mathrm i} jx} \n = {\\mathop\\sum}_{|j| \\leq C} (h, A^T e^{{\\mathrm i} j x} )_{L^2(\\mathbb T)} e^{{\\mathrm i} j x} \\,, \n$$\nwhich has the form \\eqref{forma buona resto} with $ g_j := A^T( e^{{\\mathrm i} jx } )$ and $ \\chi_j := e^{{\\mathrm i} j x} $.\nDifferentiating the second equality in \\eqref{Psiuh}, we see that \n$$\n{\\cal R}_1 (u)[h] = \\Pi_E [ A h ] \\, , \\quad A h := \\partial_u \\{ \\Psi' ( \\Pi_E u)^T \\} [\\Pi_E h, \\Pi_E \n(\\nabla H)(\\Phi(u)) ] \\,,\n$$\nwhich has the same form of $ {\\cal R}_2 $ and so \\eqref{forma buona resto}. \n\\end{proof}\n\n\n\\begin{lemma} \\label{sifulo}\nLet $ H(u ) := \\int_\\mathbb T f(u) X(u) d x $\nwhere $ X(u) = \\Pi_{E} X ( \\Pi_E u) $ and $ f(u)(x) := f( u(x)) $ is the composition operator for a function of class\n$ C^2 $. Then\n\\begin{equation}\\label{lint1}\n(\\partial_u \\nabla H) (u) [h] = f''(u) X(u) \\, h + {\\cal R} (u) [h] \n\\end{equation}\nwhere $ {\\cal R} (u) $ has the form \\eqref{forma buona resto}\nwith $ \\chi_j (u) = e^{{\\mathrm i} j x} $ or $ g_j(u) = e^{{\\mathrm i} j x} $.\n\\end{lemma}\n\n\\begin{proof}\nA direct calculation proves that $\\nabla H(u) = f'(u) X(u) + X'(u)^T [f(u)]$, \nand \\eqref{lint1} follows with\n$ {\\cal R} (u) [h] = $ $ f'(u) X'(u)[h] + $ $ \\partial_u \\{ X'(u)^T\\} [h, f(u)] + $ $ X'(u)^T [ f'(u) h ]$, \nwhich has the form \\eqref{forma buona resto}. \n\\end{proof}\n\nWe conclude this section with a technical lemma used from the end of section \\ref{step3} about the decay norms of ``finite dimensional operators\".\nNote that operators of the form \\eqref{forma buona con gli integrali} (that will appear in section \\ref{step1}) \nreduce to those in \\eqref{forma buona resto} when the functions $ g_j(\\tau) $, $ \\chi_j (\\tau)$ are independent of $ \\tau $\n\n\n\\begin{lemma}\\label{remark : decay forma buona resto}\nLet $ {\\cal R} $ be an operator of the form\n\\begin{equation}\\label{forma buona con gli integrali}\n{\\cal R} h = \\sum_{|j| \\leq C } \\int_0^1 \\big(h\\,,\\,g_j(\\tau) \\big)_{L^2(\\mathbb T)} \\chi_j (\\tau)\\,d \\tau\\,,\n\\end{equation}\nwhere the functions $g_j(\\tau),\\,\\chi_j(\\tau) \\in H^s$, $\\tau \\in [0, 1]$ depend in a Lipschitz way on the parameter $\\omega$. Then its matrix $ s$-decay norm (see \\eqref{matrix decay norm}-\\eqref{matrix decay norm Lip}) satisfies \n$$ \n| {\\cal R} |_s^{\\mathrm{Lip}(\\g)} \\leq_s {\\mathop \\sum}_{|j| \\leq C} {\\rm sup}_{\\tau \\in [0,1]} \\big\\{ \\| \\chi_j(\\tau) \\|_s^{\\mathrm{Lip}(\\g)} \\| g_j(\\tau) \\|_{s_0}^{\\mathrm{Lip}(\\g)} \n+ \\| \\chi_j(\\tau) \\|_{s_0}^{\\mathrm{Lip}(\\g)} \\| g_j(\\tau) \\|_s^{\\mathrm{Lip}(\\g)} \\big\\} \\, . \n$$ \n\\end{lemma}\n\n\\begin{proof}\nFor each $\\tau \\in [0, 1]$, the operator $ h \\mapsto (h,g_j(\\tau)) \\chi_j(\\tau) $ is the \ncomposition $ \\chi_j(\\tau) \\circ \\Pi_0 \\circ g_j(\\tau) $ of the multiplication operators for $ g_j(\\tau), \\chi_j(\\tau) $ and \n$ h \\mapsto \\Pi_0 h := \\int_{\\mathbb T} h dx $. Hence\nthe lemma follows by the interpolation estimate \\eqref{interpm Lip} \nand \\eqref{multiplication Lip}. \n\\end{proof}\n\n\\subsection{Composition with the map $G_\\d$} \\label{section:appr}\n\nIn the sequel we shall use that \n$ {\\mathfrak{I}}_\\d := {\\mathfrak{I}}_\\d (\\varphi ; \\omega) := i_\\d (\\varphi; \\omega ) - (\\varphi,0,0) $ satisfies, \nby Lemma \\ref{toro isotropico modificato} and \\eqref{ansatz 0}, \n\\begin{equation}\\label{ansatz delta}\n\\| {\\mathfrak I}_\\d \\|_{s_0+\\mu}^{{\\mathrm{Lip}(\\g)}} \n\\leq C\\varepsilon^{6 - 2b} \\gamma^{-1}\\, . \n\\end{equation}\nWe now study the Hamiltonian \n$ K := H_\\varepsilon \\circ G_\\d = \\varepsilon^{-2b} \\mathcal{H} \\circ A_\\varepsilon \\circ G_\\d $ defined in \\eqref{KHG}, \\eqref{def H eps}. \n\nRecalling \\eqref{def A eps} and \\eqref{trasformazione modificata simplettica} the map $A_\\varepsilon \\circ G_\\d$ has the form \n\\begin{equation} \\label{A eps G delta}\nA_\\varepsilon \\circ G_\\d(\\psi, \\eta, w) \n= \\varepsilon \\sum_{j \\in S} \\sqrt{\\xi_j + \\varepsilon^{2(b-1)} |j| [ y_\\d(\\psi) + L_1(\\psi) \\eta + L_2(\\psi) w ]_j } \\, e^{{\\mathrm i} [\\theta_0(\\psi)]_j} e^{{\\mathrm i} jx} + \\varepsilon^b (z_0(\\psi) + w)\n\\end{equation}\nwhere \n\\begin{equation}\\label{L1 L2}\nL_1(\\psi) := [\\partial_\\psi \\theta_0(\\psi)]^{-T} \\, , \\quad\nL_2(\\psi) := \\big[ (\\partial_\\theta \\tilde{z}_0) (\\theta_0(\\psi)) \\big]^T \\partial_x^{-1} \\, . \n\\end{equation}\n By Taylor's formula, we develop \\eqref{A eps G delta} in $w$ at $\\eta=0$, $w=0$, and we get\n$ A_\\varepsilon \\circ G_\\d(\\psi, 0, w) = $ $ T_\\delta(\\psi) + T_1(\\psi) w + T_2(\\psi)[w,w] + $ $ T_{\\geq 3}(\\psi, w) $, \nwhere \n\\begin{equation}\\label{T0}\nT_\\delta(\\psi) := (A_\\varepsilon \\circ G_\\d)(\\psi, 0, 0)\n= \\varepsilon v_\\d(\\psi) + \\varepsilon^b z_0(\\psi) \\, , \\ \\ v_\\d (\\psi):= \\sum_{j \\in S} \\sqrt{\\xi_j + \\varepsilon^{2(b-1)} |j| [ y_\\d(\\psi) ]_j } \\, e^{{\\mathrm i} [\\theta_0(\\psi)]_j} e^{{\\mathrm i} jx}\n\\end{equation}\nis the approximate isotropic torus in phase space (it corresponds to $ i_\\d $ in Lemma \\ref{toro isotropico modificato}), \n\\begin{align}\nT_1(\\psi) w & = \\varepsilon \\sum_{j \\in S} \n\\frac{\\varepsilon^{2(b-1)} |j| [ L_2(\\psi) w ]_j \\, e^{{\\mathrm i} [\\theta_0(\\psi)]_j}} \n{2 \\sqrt{ \\xi_j + \\varepsilon^{2(b-1)} |j| [ y_\\d(\\psi) ]_j }} \\, e^{{\\mathrm i} jx} \n+ \\varepsilon^b w =: \\varepsilon^{2b-1} U_1 (\\psi) w + \\varepsilon^b w \\, \\label{T1} \\\\\nT_2(\\psi)[w,w] & \n= - \\varepsilon \\sum_{j \\in S} \n\\frac{\\varepsilon^{4(b-1)} j^2 [ L_2(\\psi) w ]_j^2 \\, e^{{\\mathrm i} [\\theta_0(\\psi)]_j}} \n{8 \\{ \\xi_j + \\varepsilon^{2(b-1)} |j| [ y_\\d(\\psi) ]_j \\}^{3\/2} } \\, \ne^{{\\mathrm i} jx} =: \\varepsilon^{4b - 3} U_2(\\psi)[w,w] \\label{T2}\n\\end{align}\nand $T_{\\geq 3}(\\psi, w)$ collects all the terms of order at least cubic in $w$. In the notation of \\eqref{def A eps},\nthe function $v_\\d(\\psi) $ in \\eqref{T0} is \n$v_\\d(\\psi) = v_\\varepsilon( \\theta_0(\\psi), y_\\d(\\psi))$.\nThe terms $U_1, U_2 = O(1)$ in $\\varepsilon$.\nMoreover, using that $ L_2 (\\psi) $ in \\eqref{L1 L2} vanishes as $ z_0 = 0 $, \nthey satisfy\n\\begin{equation}\\label{extra piccolezza}\n\\| U_1 w \\|_s \\leq \\| {\\mathfrak{I}}_\\d \\|_s \\| w \\|_{s_0} + \\| {\\mathfrak{I}}_\\d \\|_{s_0} \\| w \\|_s \\, , \\quad\n\\| U_2 [w,w] \\|_s \\leq \\| {\\mathfrak{I}}_\\d \\|_s \\| {\\mathfrak{I}}_\\d \\|_{s_0} \\| w \\|_{s_0}^2 + \n\\| {\\mathfrak{I}}_\\d \\|_{s_0}^2 \\| w \\|_{s_0} \\| w \\|_s \n\\end{equation}\nand also in the $ \\| \\ \\|_s^{\\mathrm{Lip}(\\g)} $-norm.\n\nBy Taylor's formula \n$ \\mathcal{H}(u+h) \n= \\mathcal{H}(u) + ( (\\nabla \\mathcal{H})(u), h )_{L^2(\\mathbb T)} \n+ \\frac12 ( (\\partial_u \\nabla \\mathcal{H})(u) [h], h )_{L^2(\\mathbb T)} \n+ O(h^3) $. \nSpecifying at $u = T_\\delta(\\psi)$ and $ h = T_1(\\psi) w + T_2(\\psi)[w,w] + T_{\\geq 3}(\\psi,w)$, \nwe obtain that the sum of all the components of $ K = \\varepsilon^{-2b} (\\mathcal{H} \\circ A_\\varepsilon \\circ G_\\d)(\\psi, 0, w) $ \nthat are quadratic in $w$ is \n$$\n\\frac12 ( K_{02}w, w )_{L^2(\\mathbb T)} \n= \\varepsilon^{-2b} ( (\\nabla \\mathcal{H})(T_\\delta ), T_2 [w,w] )_{L^2(\\mathbb T)} \n+ \\varepsilon^{-2b} \\frac12 ( (\\partial_u \\nabla \\mathcal{H})(T_\\delta ) [T_1 w], \nT_1 w )_{L^2(\\mathbb T)} \\, . \n$$\nInserting the expressions \\eqref{T1}, \\eqref{T2} we get\n\\begin{align}\nK_{02}(\\psi) w \n& = (\\partial_u \\nabla \\mathcal{H})(T_\\delta) [w] \n+ 2 \\varepsilon^{b-1} (\\partial_u \\nabla \\mathcal{H})(T_\\delta) [U_1 w] \n+ \\varepsilon^{2(b-1)} U_1^T (\\partial_u \\nabla \\mathcal{H})(T_\\delta) [U_1 w] \\nonumber\n\\\\ & \\quad \n+ 2 \\varepsilon^{2b- 3} U_2[w, \\cdot]^T (\\nabla \\mathcal{H})(T_\\delta). \\label{K02}\n\\end{align}\n\n\\begin{lemma}\\label{dopo l'approximate inverse}\n\\begin{equation}\\label{piccolezza resti}\n ( K_{02}(\\psi) w, w )_{L^2(\\mathbb T)} \n= ( (\\partial_u \\nabla \\mathcal{H})(T_\\delta) [w], w )_{L^2(\\mathbb T)} \n+ ( R(\\psi) w, w )_{L^2(\\mathbb T)}\n\\end{equation}\nwhere $R(\\psi)w $ has the ``finite dimensional\" form \n\\begin{equation}\\label{forma buona resto con psi}\nR(\\psi) w = {\\mathop\\sum}_{|j| \\leq C} \\big( w , g_j(\\psi) \\big)_{L^2(\\mathbb T)} \\chi_j(\\psi) \n\\end{equation}\nwhere, for some $\\sigma := \\sigma (\\nu, \\tau) > 0$, \n\\begin{align}\n\\label{piccolo FBR}\n \\| g_j \\|_s^{\\mathrm{Lip}(\\g)} \\| \\chi_j \\|_{s_0}^{\\mathrm{Lip}(\\g)} + \\| g_j \\|_{s_0}^{\\mathrm{Lip}(\\g)} \\| \\chi_j \\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^{b+1} \\| {\\mathfrak I}_\\delta \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} \\\\\n\\| \\partial_i g_j [\\widehat \\imath ]\\|_s \\| \\chi_j \\|_{s_0} \n+ \\| \\partial_i g_j [\\widehat \\imath ]\\|_{s_0} \\| \\chi_j \\|_{s} + \\| g_j \\|_{s_0} \\| \\partial_i \\chi_j [\\widehat \\imath ] \\|_s \n+ \\| g_j \\|_{s} \\| \\partial_i \\chi_j [\\widehat \\imath ]\\|_{s_0} \n& \\leq_s \\varepsilon^{b + 1} \\| \\widehat \\imath \\|_{s + \\sigma}\\label{derivata piccolo FBR} \\\\ \n& + \\varepsilon^{2b-1} \\| {\\mathfrak I}_\\delta\\|_{s + \\sigma} \\|\\widehat \\imath \\|_{s_0 + \\sigma} \\,, \\nonumber\n\\end{align}\nand, as usual, $i = (\\theta, y, z)$ (see \\eqref{embedded torus i}), \n$\\widehat \\imath = (\\widehat \\theta, \\widehat y, \\widehat z)$.\n\\end{lemma}\n\n\\begin{proof}\nSince $ U_1 = \\Pi_S U_1 $ and $ U_2 = \\Pi_S U_2 $, \nthe last three terms in \\eqref{K02} have all the form \\eqref{forma buona resto con psi}\n(argue as in Lemma \\ref{lemma astratto potente}). \nWe now prove that they are also small in size. \n\nThe contributions in \\eqref{K02} from $ H_2 $ are better analyzed by the expression\n$$\n\\varepsilon^{-2b} H_2 \\circ A_\\varepsilon \\circ G_\\delta (\\psi, \\eta, w) = const + \\sum_{j \\in S^+} j^3 \n\\big[ y_\\delta (\\psi) + L_1(\\psi)\\eta + L_2 (\\psi) w \\big]_j + \\frac{1}{2} \\int_{\\mathbb T} ( z_0 (\\psi) + w)_x^2 \\, dx \n$$\nwhich follows by \\eqref{shape H2}, \\eqref{trasformazione modificata simplettica}, \\eqref{L1 L2}.\n Hence the only contribution to \n$ (K_{02} w,w) $ is $ \\int_{\\mathbb T} w_x^2 \\, dx $. \nNow we consider the cubic term $ \\mathcal{H}_3 $ in \\eqref{H3tilde}.\nA direct calculation shows that for $ u = v + z $, \n$ \\nabla \\mathcal{H}_3 ( u ) = 3 z^2 + 6 \\Pi_S^\\bot (v z) $, and\n$ \\partial_u \\nabla \\mathcal{H}_3 ( u ) [U_1 w] = 6 \\Pi_S^\\bot ( z U_1 w) $ (since $ U_1 w \\in H_S $).\nTherefore\n\\begin{equation}\\label{mH 3 T delta}\n\\nabla \\mathcal{H}_3(T_\\delta) = 3 \\varepsilon^{2 b} z_0^2 + 6 \\varepsilon^{b + 1} \\Pi_S^\\bot(v_\\delta z_0 ) \n\\,,\\quad \\partial_u \\nabla \\mathcal{H}_3 ( T_\\delta ) [U_1 w] = 6 \\varepsilon^b \\Pi_S^\\bot ( z_0 \\, U_1 w) \\, . \n\\end{equation}\nBy \\eqref{mH 3 T delta} one has \n$ ( (\\partial_u \\nabla \\mathcal{H}_3)(T_\\delta) [U_1 w ], U_1 w )_{L^2(\\mathbb T)} = 0 $, \nand since also $U_2 = \\Pi_S U_2$, \n\\begin{equation}\\label{contributi H3}\n\\varepsilon^{b - 1}\\partial_u \\nabla {\\cal H}_3(T_\\delta)[U_1 w] + \\varepsilon^{2 b - 3} U_2[w, \\cdot]^T \\nabla {\\cal H}_3(T_\\delta) = 6 \\varepsilon^{2 b - 1} \\Pi_S^\\bot(z_0 U_1 w) + 3 \\varepsilon^{4 b - 3} U_2[w, \\cdot]^T z_0^2 \\,.\n\\end{equation}\nThese terms have the form \\eqref{forma buona resto con psi} and, using \\eqref{extra piccolezza}, \\eqref{ansatz 0}, \nthey satisfy \\eqref{piccolo FBR}. \n\nFinally we consider all the terms which arise from ${\\cal H}_{\\geq 4} = O(u^4)$. \nThe operators\n$ \\varepsilon^{b - 1} \\partial_u \\nabla {\\cal H}_{\\geq 4}(T_\\delta) U_1 $, \n$ \\varepsilon^{2(b - 1)} U_1^T (\\partial_u \\nabla {\\cal H}_{\\geq 4})(T_\\delta) U_1$, \n$ \\varepsilon^{2 b - 3} U_2^T \\nabla {\\cal H}_{\\geq 4}(T_\\delta) $\nhave the form \\eqref{forma buona resto con psi} and, using $ \\| T_\\delta \\|_s^{\\mathrm{Lip}(\\g)} \\leq \\varepsilon (1 + \\| {\\mathfrak{I}}_\\d \\|_s^{\\mathrm{Lip}(\\g)}) $,\n\\eqref{extra piccolezza}, \\eqref{ansatz 0}, \nthe bound \\eqref{piccolo FBR} holds. Notice that the biggest term is $ \\varepsilon^{b - 1} \\partial_u \\nabla {\\cal H}_{\\geq 4}(T_\\delta) U_1 $. \n\nBy \\eqref{derivata i delta} and using explicit formulae \\eqref{L1 L2}-\\eqref{T2} we get estimate \\eqref{derivata piccolo FBR}.\n\\end{proof}\n\nThe conclusion of this section is that, after the composition with \nthe action-angle variables, the rescaling \\eqref{rescaling kdv quadratica}, \nand the transformation $ G_\\delta $, the linearized operator \nto analyze is $ H_S^\\bot \\ni w \\mapsto (\\partial_u \\nabla \\mathcal{H})(T_\\delta) [w] $, up to finite dimensional operators which have the form \\eqref{forma buona resto con psi} and size \\eqref{piccolo FBR}. \n\n\n\n\\subsection{The linearized operator in the normal directions}\n\nIn view of \\eqref{piccolezza resti} we now compute \n$ ( (\\partial_u \\nabla \\mathcal{H})(T_\\delta) [w], w )_{L^2(\\mathbb T)} $, $ w \\in H_S^\\bot $, where\n$ \\mathcal{H} = H \\circ \\Phi_B $ and $\\Phi_B $ is the Birkhoff map of Proposition \\ref{prop:weak BNF}. \nIt is convenient to estimate separately the terms in \n\\begin{equation}\\label{mH H2H3H5}\n\\mathcal{H} = H \\circ \\Phi_B = (H_2 + H_3) \\circ \\Phi_B + H_{\\geq 5} \\circ \\Phi_B \n\\end{equation}\nwhere $ H_2, H_3, H_{\\geq 5}$ are defined in \\eqref{H iniziale KdV}.\n\nWe first consider $ H_{\\geq 5} \\circ \\Phi_B $. \nBy \\eqref{H iniziale KdV} we get\n$ \\nabla H_{\\geq 5}(u) = \\pi_0[ (\\partial_u f)(x, u, u_x) ] - \\partial_x \\{ (\\partial_{u_x} f)(x, u,u_x) \\} $, see \\eqref{def pi 0}. \nSince the Birkhoff transformation $ \\Phi_B $ has the form \\eqref{finito finito}, \nLemma \\ref{lemma astratto potente} (at $ u = T_\\delta $, see \\eqref{T0}) implies that\n\\begin{align} \n\\partial_u \\nabla ( H_{\\geq 5} \\circ \\Phi_B ) (T_\\delta) [h] \n& = \n(\\partial_u \\nabla H_{\\geq 5})(\\Phi_B(T_\\delta)) [h] + {\\cal R}_{H_{\\geq 5}}(T_\\delta)[h] \n\\notag \\\\ \n& = \\partial_x (r_1(T_\\delta) \\partial_x h ) + r_0(T_\\delta) h + {\\cal R}_{H_{\\geq 5}}(T_\\delta)[h] \n\\label{der grad struttura separata5}\n\\end{align}\nwhere the multiplicative functions $r_0(T_\\d)$, $r_1(T_\\d)$ are\n\\begin{alignat}{2} \\label{r0r1 def}\nr_0 (T_\\delta) & := \\sigma_0(\\Phi_B(T_\\delta)), \\qquad & \n\\sigma_0(u) & := (\\partial_{uu} f)(x, u, u_x) - \\partial_x \\{ (\\partial_{u u_x} f)(x, u, u_x) \\}, \n\\\\\n\\label{sigma0sigma1 def}\nr_1 (T_\\delta) & := \\sigma_1(\\Phi_B(T_\\delta)), \\quad &\n\\sigma_1(u) & := - (\\partial_{u_x u_x} f)(x, u, u_x) , \n\\end{alignat}\nthe remainder $ {\\cal R}_{H_{\\geq 5}}(u) $ has the form \\eqref{forma buona resto}\nwith $\\chi_j = e^{{\\mathrm i} jx}$ or $g_j = e^{{\\mathrm i} jx}$ and, using \\eqref{resti012}, \nit satisfies, \nfor some $ \\sigma := \\sigma (\\nu, \\tau) > 0$, \n\\begin{align*}\n \\| g_j \\|_s^{\\mathrm{Lip}(\\g)} \\| \\chi_j \\|_{s_0}^{\\mathrm{Lip}(\\g)} + \\| g_j \\|_{s_0}^{\\mathrm{Lip}(\\g)} \\| \\chi_j \\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^4 (1 + \\| {\\mathfrak{I}}_\\d \\|_{s+2}^{\\mathrm{Lip}(\\g)} ) \\\\\n\\| \\partial_i g_j [\\widehat \\imath ]\\|_s \\| \\chi_j \\|_{s_0} \n+ \\| \\partial_i g_j [\\widehat \\imath ]\\|_{s_0} \\| \\chi_j \\|_{s} \n+ \\| g_j \\|_{s_0} \\| \\partial_i \\chi_j [\\widehat \\imath ] \\|_s \n+ \\| g_j \\|_{s} \\| \\partial_i \\chi_j [\\widehat \\imath ]\\|_{s_0} \n& \\leq_s \\varepsilon^4 ( \\| \\widehat \\imath \\|_{s+\\sigma} \n+ \\| {\\mathfrak{I}}_\\d \\|_{s+2} \\| \\widehat \\imath \\|_{s_0 + 2} ). \n\\end{align*}\nNow we consider the contributions from $ (H_2 + H_3) \\circ \\Phi_B $.\nBy Lemma \\ref{lemma astratto potente} and the expressions of $ H_2, H_3 $ in \\eqref{H iniziale KdV} we deduce that\n$$\n\\partial_u \\nabla ( H_2 \\circ \\Phi_B) (T_\\delta) [h] \n= - \\partial_{xx} h + {\\cal R}_{H_2}(T_\\delta)[h] \\,, \\quad \n\\partial_u \\nabla ( H_3 \\circ \\Phi_B) (T_\\delta) [h] \n= 6 \\Phi_B (T_\\delta) h + {\\cal R}_{H_3}(T_\\delta)[h] \\,,\n$$\nwhere $ \\Phi_B (T_\\delta) $ is a function with zero space average, \nbecause $ \\Phi_B: H^1_0 (\\mathbb T_x) \\to H^1_0 (\\mathbb T_x)$ (Proposition \\ref{prop:weak BNF}) \nand $ {\\cal R}_{H_2}(u) $, $ {\\cal R}_{H_3}(u) $ have the form \\eqref{forma buona resto}. \nBy \\eqref{resti012}, the size $ ( {\\cal R}_{H_2} + {\\cal R}_{H_3}) (T_\\delta ) = O( \\varepsilon ) $. \nWe expand\n$$\n( {\\cal R}_{H_2} + {\\cal R}_{H_3}) (T_\\delta ) = \\varepsilon {\\cal R}_1 + \\varepsilon^2 {\\cal R}_2 + {\\tilde {\\cal R}}_{> 2} \\,, \n$$ \nwhere $\\tilde \\mathcal{R}_{>2}$ has size $o(\\varepsilon^2)$, \nand we get, $ \\forall h \\in H_S^\\bot $, \n\\begin{equation}\\label{useful repr}\n\\Pi_S^\\bot \\partial_u \\nabla ((H_2 + H_3) \\circ \\Phi_B) (T_\\delta) [h] \n= - \\partial_{xx} h + \\Pi_S^\\bot (6 \\Phi_B (T_\\delta) h ) + \n\\Pi_S^\\bot ( \\varepsilon {\\cal R}_1 + \\varepsilon^2 {\\cal R}_2 + {\\tilde {\\cal R}}_{> 2} ) [h] \\, . \n\\end{equation}\nWe also develop the function $ \\Phi_B (T_\\delta) $ is powers of $ \\varepsilon $. \nExpand $\\Phi_B (u) = u + \\Psi_2 (u) + \\Psi_{\\geq 3}(u) $, \nwhere $ \\Psi_2 (u) $ is quadratic, $ \\Psi_{\\geq 3} (u) = O(u^3)$, and both map \n$ H_0^1(\\mathbb T_x) \\to H_0^1(\\mathbb T_x) $.\nAt $ u = T_\\d = \\varepsilon v_\\d + \\varepsilon^b z_0 $ \nwe get\n\\begin{align}\n\\Phi_B ( T_\\delta ) & = T_\\delta + \\Psi_2 (T_\\delta) + \\Psi_{\\geq 3}(T_\\delta) \n= \\varepsilon v_\\d + \\varepsilon^2 \\Psi_2 ( v_\\d ) + \\tilde q \n\\label{funzione moltiplicativa}\n\\end{align}\nwhere $ \\tilde q := \\varepsilon^b z_0 + \\Psi_2 ( T_\\d ) - \\varepsilon^2 \\Psi_2 (v_\\d) \n+ \\Psi_{\\geq 3} (T_\\d) $ \nhas zero space average and it satisfies \n$$\n\\| \\tilde q \\|_s^{\\mathrm{Lip}(\\g)} \n\\leq_s \\varepsilon^3 + \\varepsilon^b \\| {\\mathfrak{I}}_\\delta \\|_s^{\\mathrm{Lip}(\\g)}\\,,\\quad \n\\| \\partial_i \\tilde q [\\widehat \\imath ] \\|_s \n\\leq_s \\varepsilon^b \\big( \\| \\widehat \\imath \\|_s + \\| {\\mathfrak I}_\\delta \\|_s \\| \\widehat \\imath \\|_{s_0} \\big)\\,.\n$$\nIn particular, its low norm $\\| \\tilde q \\|_{s_0}^{\\mathrm{Lip}(\\g)} \\leq_{s_0} \\varepsilon^{6-b} \\gamma^{-1} = o(\\varepsilon^2)$. \n\nWe need an exact expression of the terms of order $ \\varepsilon $ and $ \\varepsilon^2 $ in \\eqref{useful repr}. \nWe compare the Hamiltonian \\eqref{widetilde cal H} with \\eqref{mH H2H3H5}, \nnoting that $ (H_{\\geq 5} \\circ \\Phi_B)(u) = O(u^5) $ \nbecause $ f $ satisfies \\eqref{order5} and $ \\Phi_B (u) = O(u) $. Therefore \n$$\n(H_2 + H_3) \\circ \\Phi_B = H_2 + \\mathcal{H}_3 + \\mathcal{H}_4 + O(u^5) \\,,\n$$\nand the homogeneous terms of $ (H_2 + H_3) \\circ \\Phi_B $ of degree $ 2, 3, 4 $ in $ u $ are $ H_2 $, $\\mathcal{H}_3 $, $ \\mathcal{H}_4 $ respectively. \nAs a consequence, the terms of order $ \\varepsilon $ and $ \\varepsilon^2 $ in \\eqref{useful repr} (both in the function \n$ \\Phi_B (T_\\d ) $ and in the remainders $ \\mathcal{R}_1, \\mathcal{R}_2 $) come only from $ H_2 + \\mathcal{H}_3 + \\mathcal{H}_4 $. \nActually they come from $ H_2 $, $ \\mathcal{H}_3 $ and $ \\mathcal{H}_{4,2} $ (see \\eqref{H3tilde}, \\eqref{mH3 mH4})\nbecause, at $ u = T_\\d = \\varepsilon v_\\d + \\varepsilon^b z_0 $, \nfor all $ h \\in H_S^\\bot $, \n$$\n\\Pi_S^\\bot (\\partial_u \\nabla \\mathcal{H}_4)(T_\\delta) [h] = \n\\Pi_S^\\bot (\\partial_u \\nabla \\mathcal{H}_{4,2})(T_\\delta) [h] + o(\\varepsilon^2) \\,. \n$$\nA direct calculation based on the expressions \\eqref{H3tilde}, \\eqref{mH3 mH4}\nshows that, for all $ h \\in H_S^\\bot $, \n\\begin{align}\n\\Pi_S^\\bot (\\partial_u \\nabla ( H_2 + \\mathcal{H}_3 + \\mathcal{H}_4)) (T_\\delta)[h] \n& = - \\partial_{xx} h + 6 \\varepsilon \\Pi_S^\\bot ( v_\\d h ) + 6 \\varepsilon^b \\Pi_S^\\bot ( z_0 h ) \n+ \\varepsilon^2 \\Pi_S^\\bot \\big\\{ 6 \\pi_0 [ (\\partial_x^{-1} v_\\d)^2 ] h \\nonumber \\\\\n& \n+ 6 v_\\d \\Pi_S [ (\\partial_x^{-1} v_\\d) (\\partial_x^{-1} h) ] -\n 6 \\partial_x^{-1} \\{ (\\partial_x^{-1} v_\\d) \\Pi_S [v_\\d h] \\} \\big\\} + \no(\\varepsilon^2). \\label{lin basso}\n\\end{align}\nThus, comparing the terms of order $ \\varepsilon, \\varepsilon^2 $ \nin \\eqref{useful repr} (using \\eqref{funzione moltiplicativa}) with those in \\eqref{lin basso} we deduce that the operators $\\mathcal{R}_1, \\mathcal{R}_2$ and the function $\\Psi_2 (v_\\d)$ are\n\\begin{equation}\\label{R nullo1R2}\n{\\cal R}_1 = 0 , \\quad \n{\\cal R}_2 [h ] = 6 v_\\d \\Pi_S \\big[ (\\partial_x^{-1} v_\\d) (\\partial_x^{-1} h) \\big] \n- 6 \\partial_x^{-1} \\{ (\\partial_x^{-1} v_\\d) \\Pi_S [v_\\d h] \\} \\,, \\quad \n\\Psi_2 (v_\\d) = \\pi_0 [ (\\partial_x^{-1} v_\\d)^2 ] . \n\\end{equation}\nIn conclusion, by \\eqref{mH H2H3H5}, \n\\eqref{useful repr}, \\eqref{der grad struttura separata5},\n\\eqref{funzione moltiplicativa}, \\eqref{R nullo1R2}, we get, \nfor all $ h \\in H_{S^\\bot } $, \n\\begin{align}\n\\Pi_S^\\bot \\partial_u \\nabla \\mathcal{H} (T_\\delta)[h] \n& = - \\partial_{xx} h + \\Pi_S^\\bot \n\\big[ \\big( \\varepsilon 6 v_\\d + \n\\varepsilon^2 6 \\pi_0 [ (\\partial_x^{-1} v_\\d)^2 ] + q_{>2} + p_{\\geq 4 } \n\\big) h \\big] \\nonumber \\\\\n& \\quad + \n \\Pi_S^\\bot \\partial_x (r_1(T_\\delta) \\partial_x h ) + \n \\varepsilon^2 \\Pi_S^\\bot {\\cal R}_2[h] + \\Pi_S^\\bot {\\cal R}_{> 2} [h] \\label{lafinalina}\n\\end{align}\nwhere $ r_1 $ is defined in \\eqref{r0r1 def}, $ {\\cal R}_2 $ in \\eqref{R nullo1R2}, the remainder \n$ {\\cal R}_{> 2} := {\\tilde {\\cal R}}_{> 2} + {\\cal R}_{H_{\\geq 5}} (T_\\d) $ \nand the functions\n(using also \\eqref{r0r1 def}, \\eqref{sigma0sigma1 def}, \\eqref{order5}), \n\\begin{align}\\label{def p3}\nq_{>2} & := 6 \\tilde q \n + \\varepsilon^3 \\big( (\\partial_{uu} f_5)(v_\\d, (v_\\d)_x) - \\partial_x \\{ (\\partial_{u u_x} f_5)(v_\\d, (v_\\d)_x) \\} \\big) \\\\\np_{\\geq 4 } & := \nr_0 (T_\\d) - \\varepsilon^3 \\big[ (\\partial_{uu} f_5)(v_\\d, (v_\\d)_x) - \\partial_x \\{ (\\partial_{u u_x} f_5)(v_\\d, (v_\\d)_x) \\} \\big] \\, . \n \\label{def p>4} \n\\end{align}\n\n\\begin{lemma}\\label{p3 zero average}\n$ \\int_{\\mathbb T} q_{>2} dx = 0 $. \n\\end{lemma}\n\n\\begin{proof}\nWe already\nobserved that $ \\tilde q $ has zero $x$-average as well as\nthe derivative $ \\partial_x \\{ (\\partial_{u u_x} f_5)(v, v_x) \\} $.\nFinally \n\\begin{equation}\\label{unico pezzo}\n(\\partial_{uu} f_5)(v, v_x) = \\sum_{j_1, j_2, j_3 \\in S } c_{j_1j_2j_3} v_{j_1} v_{j_2} v_{j_3} e^{{\\mathrm i} (j_1+ j_2 + j_3) x } \\, , \n\\quad v := \\sum_{j \\in S} v_j e^{{\\mathrm i} j x}\n\\end{equation}\nfor some coefficient $ c_{j_1j_2j_3} $, and therefore it\nhas zero average by hypothesis (${\\mathtt S}1$). \n\\end{proof}\n\nBy Lemma \\ref{dopo l'approximate inverse}\nand the results of this section (in particular \\eqref{lafinalina}) we deduce: \n\n\\begin{proposition}\\label{prop:lin}\nAssume \\eqref{ansatz delta}. Then the Hamiltonian operator $ {\\cal L}_\\omega $ has the form, \n$ \\forall h \\in H_{S^\\bot}^s ( \\mathbb T^{\\nu+1}) $,\n\\begin{equation}\\label{Lom KdVnew}\n{\\cal L}_\\omega h := \\omega \\!\\cdot \\!\\partial_{\\varphi} h - \\partial_x K_{02} h \n= \\Pi_S^\\bot \\big( \\omega \\! \\cdot \\! \\partial_\\varphi h + \\partial_{xx} (a_1 \\partial_x h) +\n\\partial_x ( a_0 h ) \n- \\varepsilon^2 \\partial_x {\\cal R}_2 h - \\partial_x \\mathcal{R}_* h \\big) \n\\end{equation}\nwhere $ {\\cal R}_2 $ is defined in \\eqref{R nullo1R2}, \n$ {\\mathcal{R}}_* := {\\cal R}_{> 2} + R(\\psi) $ (with $R(\\psi)$ defined in Lemma \\ref{dopo l'approximate inverse}), the functions \n\\begin{equation}\\label{a1p1p2}\na_1 := 1 - r_1 ( T_\\delta ) \\, , \\quad \na_0 := - ( \\varepsilon p_1 + \\varepsilon^2 p_2 + q_{>2} + p_{\\geq 4} ) \\, , \\quad \np_1 := 6 v_\\d \\, , \\quad p_2 := 6 \\pi_0 [ (\\partial_x^{-1} v_\\d)^2 ]\\,, \n\\end{equation}\nthe function $ q_{>2} $ is defined in \\eqref{def p3} and satisfies $ \\int_{\\mathbb T} q_{>2} dx = 0 $, the function $ p_{ \\geq 4} $ is defined in \\eqref{def p>4}, \n$ r_1 $ in \\eqref{sigma0sigma1 def}, \n$ T_\\delta $ and $ v_\\d $ in \\eqref{T0}.\nFor $ p_k = p_1, p_2 $, \n\\begin{alignat}{2} \\label{stime pk}\n\\| p_k \\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s 1 + \\| {\\mathfrak I}_\\d \\|_s^{\\mathrm{Lip}(\\g)}, \n& \\quad \\qquad \n\\| \\partial_i p_k [ \\widehat \\imath ] \\|_s \n& \\leq_s \\| \\widehat \\imath \\|_{s+1} + \\| {\\mathfrak I}_\\delta \\|_{s+1} \\| \\widehat \\imath \\|_{s_0+1}, \n\\\\\n\\label{stima q>2}\n\\| q_{>2} \\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^3 + \\varepsilon^b \\| {\\mathfrak{I}}_\\delta \\|_{s}^{\\mathrm{Lip}(\\g)}\\,,\n& \\quad \n\\| \\partial_i q_{>2} [\\widehat \\imath ] \\|_s \n& \\leq_s \\varepsilon^b \\big( \\| \\widehat \\imath \\|_{s+1} + \\| {\\mathfrak I}_\\delta \\|_{s+1} \\| \\widehat \\imath \\|_{s_0+1} \\big),\n\\\\\n\\| a_1 -1 \\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^3 \\big( 1 + \\| {\\mathfrak I}_\\d \\|_{s+1}^{\\mathrm{Lip}(\\g)} \\big) \\,, \n& \\quad \n\\| \\partial_i a_1[\\widehat \\imath ] \\|_s \n& \\leq_s \\varepsilon^3 \\big( \\| \\widehat \\imath \\|_{s + 1} + \\| {\\mathfrak I}_\\delta \\|_{s + 1} \\| \\widehat \\imath \\|_{s_0 + 1} \\big) \\label{stima a1} \\\\\n\\| p_{\\geq 4} \\|_s^{\\mathrm{Lip}(\\g)} & \\leq_s \\varepsilon^4 + \\varepsilon^{b + 2} \\| {\\mathfrak{I}}_\\delta \\|_{s + 2}^{\\mathrm{Lip}(\\g)}\\,,& \\quad \\| \n\\partial_i p_{\\geq 4}[\\widehat \\imath ] \\|_s & \\leq_s \\varepsilon^{b + 2} \\big( \\| \\widehat \\imath \\|_{s + 2} + \\| {\\mathfrak{I}}_\\delta\\|_{s + 2} \\| \\widehat \\imath \\|_{s_0 + 2} \\big) \\label{p geq 4}\n\\end{alignat}\nwhere $ {\\mathfrak{I}}_\\d (\\varphi) := (\\theta_0(\\varphi) - \\varphi, y_\\d(\\varphi), z_0(\\varphi)) $ corresponds to $T_\\delta$. \nThe remainder $ {\\cal R}_{2} $ has the form \\eqref{forma buona resto} with \n\\begin{align} \\label{stime resto 1}\n\\| g_j \\|_s^{\\mathrm{Lip}(\\g)} + \\| \\chi_j \\|_s^{\\mathrm{Lip}(\\g)} \n\\leq_s 1+ \\| {\\mathfrak I}_\\delta \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} \\, , \\quad \n\\| \\partial_i g_j [\\widehat \\imath ]\\|_s + \\| \\partial_i \\chi_j [\\widehat \\imath ]\\|_s\n \\leq_s \\| \\widehat \\imath \\|_{s + \\sigma} + \\| {\\mathfrak I}_\\delta\\|_{s + \\sigma} \\|\\widehat \\imath \\|_{s_0 + \\sigma} \n\\end{align}\nand also $ {\\cal R}_* $ has the form \\eqref{forma buona resto} with\n\\begin{align} \\label{stima cal R*}\n \\| g_j^* \\|_s^{\\mathrm{Lip}(\\g)} \\| \\chi_j^* \\|_{s_0}^{\\mathrm{Lip}(\\g)} + \\| g_j^* \\|_{s_0}^{\\mathrm{Lip}(\\g)} \\| \\chi_j^* \\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^3 + \\varepsilon^{b+1} \\| {\\mathfrak I}_\\delta \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)}\n \\\\\n\\| \\partial_i g_j^* [\\widehat \\imath ]\\|_s \\| \\chi_j^* \\|_{s_0} \n+ \\| \\partial_i g_j^* [\\widehat \\imath ]\\|_{s_0} \\| \\chi_j^* \\|_{s} + \\| g_j^* \\|_{s_0} \\| \\partial_i \\chi_j^* [\\widehat \\imath ] \\|_s \n+ \\| g_j^* \\|_{s} \\| \\partial_i \\chi_j^* [\\widehat \\imath ]\\|_{s_0} \n& \\leq_s \\varepsilon^{b + 1} \\| \\widehat \\imath \\|_{s + \\sigma}\n\\label{derivate stima cal R*} \\\\ \n& + \\varepsilon^{2b-1} \\| {\\mathfrak I}_\\delta\\|_{s + \\sigma} \\|\\widehat \\imath \\|_{s_0 + \\sigma} \\nonumber \\, . \n\\end{align} \n\\end{proposition}\n\nThe bounds \\eqref{stime resto 1}, \\eqref{stima cal R*} imply, by \nLemma \\ref{remark : decay forma buona resto}, estimates \nfor the $ s $-decay norms of $ {\\cal R}_2 $ and \n$ {\\cal R}_* $. The linearized operator $ {\\cal L}_\\omega := {\\cal L}_\\omega (\\omega, i_\\d (\\omega))$ \ndepends on the parameter $ \\omega $ both directly and also through the dependence on the \ntorus $ i_\\d (\\omega ) $.\nWe have estimated also the partial derivative $ \\partial_i $ with respect to the variables $ i $ (see \\eqref{embedded torus i})\nin order to \ncontrol, along the nonlinear Nash-Moser iteration, \nthe Lipschitz variation of the eigenvalues of $ {\\cal L}_\\omega $ with respect to $ \\omega $\nand the approximate solution $ i_\\d $. \n\n\n\n\\section{Reduction of the linearized operator in the normal directions}\\label{operatore linearizzato sui siti normali}\n\nThe goal of this section is to conjugate the Hamiltonian operator $ {\\cal L}_\\omega $ in \\eqref{Lom KdVnew} to the \ndiagonal operator $ {\\cal L}_\\infty $ defined in \\eqref{Lfinale}. \nThe proof is obtained applying different kind of symplectic transformations. We shall always assume \\eqref{ansatz delta}. \n\n\n\\subsection{Change of the space variable }\\label{step1}\n\nThe first task is to conjugate $ {\\cal L}_\\omega $ in \\eqref{Lom KdVnew} to $ {\\cal L}_1 $ in\n \\eqref{cal L1 Kdv}, which has \nthe coefficient of $\\partial_{xxx}$ independent on the space variable. \nWe look for a $ \\varphi $-dependent family of {\\it symplectic} diffeomorphisms $\\Phi (\\varphi) $ of $ H_S^\\bot $ \nwhich differ from \n\\begin{equation}\\label{primo cambio di variabile modi normali}\n{\\cal A}_{\\bot} := \\Pi_S^\\bot {\\cal A} \\Pi_S^\\bot \\, , \\quad \n({\\cal A} h)(\\varphi,x) := (1 + \\beta_x(\\varphi,x)) h(\\varphi,x + \\beta(\\varphi,x)) \\, , \n\\end{equation}\nup to a small ``finite dimensional\" remainder, see \\eqref{forma buona resto cambio di variabile hamiltoniano}.\nEach \n$ {\\cal A}(\\varphi) $ is a symplectic map\nof the phase space, see \\cite{BBM}-Remark 3.3. \nIf $ \\| \\b \\|_{W^{1,\\infty}} < 1 \/ 2 $ then \n $ {\\cal A} $ is invertible, see Lemma \\ref{lemma:utile}, and \nits inverse and adjoint maps are \n\\begin{equation}\\label{cambio di variabile inverso}\n({\\cal A}^{-1} h)(\\varphi,y) := (1 + \\tilde{\\beta}_y(\\varphi,y)) h(\\varphi, y + \\tilde{\\beta}(\\varphi,y)) \\,, \n\\quad \n({\\cal A}^T h) (\\varphi,y) = h(\\varphi, y + \\tilde{\\beta}(\\varphi,y)) \n\\end{equation}\nwhere\n$ x = y + \\tilde{\\b} (\\varphi, y) $ is the inverse diffeomorphism (of $\\mathbb T$) of $ y = x + \\b (\\varphi, x) $. \n\nThe restricted maps $ {\\cal A}_\\bot (\\varphi): H_S^\\bot \\to H_S^\\bot $ are not symplectic. \nIn order to find a symplectic diffeomorphism near $ {\\cal A}_\\bot (\\varphi) $, \nthe first observation is \nthat each $ {\\cal A }(\\varphi ) $ can be seen as the time $1$-flow of a time dependent \nHamiltonian PDE. Indeed $ {\\cal A }(\\varphi ) $ (for simplicity we skip the dependence on $ \\varphi $) \nis homotopic to the identity \nvia the path of symplectic diffeomorphisms \n$$ \nu \\mapsto \n(1+ \\tau \\b_x ) u ( x+ \\tau \\b(x) ), \\quad \\tau \\in [0,1 ] \\, , \n$$\nwhich is the trajectory solution of the time dependent, linear Hamiltonian PDE\n\\begin{equation} \\label{transport-free}\n\\partial_\\tau u = \\partial_x (b(\\tau, x) u) \\, , \\quad b (\\tau, x) := \\frac{\\beta(x)}{1 + \\tau \\beta_x(x)}\\, , \n\\end{equation}\nwith value $ u (x) $ at $ \\tau = 0 $ and $ {\\cal A}u = (1+ \\b_x(x) ) u ( x+ \\b(x) ) $ at $ \\t = 1 $. \nThe equation \\eqref{transport-free} is a {\\it transport} equation. Its \nassociated\ncharactheristic ODE is \n \\begin{equation}\\label{equazione delle caratteristiche}\n \\frac{d}{d\\tau} x = - b(\\tau, x ) \\, .\n \\end{equation}\nWe denote its flow by $ \\gamma^{\\tau_0, \\t} $, namely $ \\gamma^{\\tau_0, \\t} (y) $ is the \n solution of \\eqref{equazione delle caratteristiche} with $ \\gamma^{\\tau_0, \\tau_0} (y) = y $. \nEach $ \\gamma^{\\tau_0, \\t} $ is a diffeomorphism of the torus $ \\mathbb T_x $. \n\n\\begin{remark}\\label{rem:ca}\nLet $ y \\mapsto y + \\tilde \\beta(\\tau, y)$ be the inverse diffeomorpshim of \n$x \\mapsto x + \\tau \\beta(x)$. \nDifferentiating the identity $\\tilde \\beta(\\tau, y) + \\tau \\beta(y + \\tilde \\beta(\\tau, y)) = 0$ with respect to $\\tau$ it results that \n $ \\gamma^\\t (y) := \\gamma^{0,\\tau} (y) = y + \\tilde \\beta(\\tau, y) $. \n\\end{remark}\n\nThen we define a symplectic map $\\Phi $ of $ H_S^\\bot $ as the time-1 flow of the Hamiltonian PDE\n \\begin{equation}\\label{problemi di cauchy} \n \\partial_\\tau u = \\Pi_S^\\bot \\partial_x (b(\\tau, x) u) = \\partial_x (b(\\tau, x) u) - \\Pi_S \\partial_x (b(\\tau, x) u) \\, ,\n \\quad u \\in H_S^\\bot \\, .\n \\end{equation}\nNote that $ \\Pi_S^\\bot \\partial_x (b(\\tau, x) u) $ is the Hamiltonian vector field generated by \n$ \\frac12 \\int_{\\mathbb T} b(\\tau, x) u^2 dx $ restricted to $ H_S^\\bot $. \nWe denote by $ \\Phi^{\\tau_0,\\tau} $ the flow of \n \\eqref{problemi di cauchy}, namely $ \\Phi^{\\tau_0, \\tau} (u_0 ) $ is the solution of \\eqref{problemi di cauchy} \n with initial condition $ \\Phi^{\\tau_0, \\tau_0} (u_0 ) = u_0 $.\nThe flow is well defined in Sobolev spaces \n$ H^s_{S^\\bot} (\\mathbb T_x) $ for $ b(\\tau, x) $ is smooth enough\n(standard theory of linear hyperbolic PDEs, see e.g. \nsection 0.8 in \\cite{Taylor}). \n It is natural to expect that the \ndifference between the flow map $ \\Phi := \\Phi^{0,1} $ and $ {\\cal A}_\\bot $ is a ``finite-dimensional\" remainder\nof the size of $ \\b $. \n\n\\begin{lemma}\\label{modifica simplettica cambio di variabile}\nFor $ \\| \\beta \\|_{W^{s_0 + 1,\\infty}} $ small, there exists\nan invertible symplectic transformation $ \\Phi = {\\cal A}_\\bot + {\\cal R}_\\Phi $ of $ H_{S^\\bot}^s $, \n where $ {\\cal A}_\\bot $ is defined in \\eqref{primo cambio di variabile modi normali} \nand $ {\\cal R}_\\Phi $ is a ``finite-dimensional\" remainder \n\\begin{equation}\\label{forma buona resto cambio di variabile hamiltoniano}\n{\\cal R}_\\Phi h= \\sum_{j \\in S} \\int_0^1 (h, g_j (\\tau) )_{L^2(\\mathbb T)} \\chi_j (\\tau) d \\tau + \\sum_{j \\in S} \\big(h, \\psi_j \\big)_{L^2(\\mathbb T)} e^{{\\mathrm i} j x}\n\\end{equation}\nfor some functions $ \\chi_j (\\tau), g_j (\\tau) , \\psi_j \\in H^s $ satisfying \n\\begin{equation}\\label{stime forma buona resto cambio di variabile hamiltoniano}\n\\| \\psi_j\\|_s\\,,\\, \\| g_j(\\tau)\\|_s \\leq_s \\| \\beta\\|_{W^{s + 2, \\infty}}\\,,\n\\quad \\| \\chi_j(\\tau)\\|_s \\leq_s 1 + \\| \\beta \\|_{W^{s + 1, \\infty}} \\,,\\quad \\forall \\tau \\in [0, 1]\\, .\n\\end{equation}\nFurthermore, \nthe following tame estimates holds\n\\begin{equation}\\label{stime Phi Phi -1}\n\\| \\Phi^{\\pm 1}h\\|_s \\leq_s \\| h \\|_s + \\| \\beta \\|_{W^{s + 2, \\infty}} \\| h \\|_{s_0} \\, , \\quad \\forall h \\in H^s_{S^\\bot} \\, . \n\\end{equation}\n\\end{lemma}\n\n\\begin{proof}\nLet $ w (\\tau, x ) := (\\Phi^\\t u_0)(x) $ denote the solution of \\eqref{problemi di cauchy} with initial condition \n$ \\Phi^0 (w) = u_0 \\in H_S^\\bot $. \nThe difference \n\\begin{equation}\\label{differenza voluta}\n({\\cal A}_\\bot - \\Phi) u_0 = \\Pi_S^\\bot {\\cal A} u_0 - w(1, \\cdot) \n= {\\cal A}u_0 - w(1, \\cdot ) - \\Pi_S {\\cal A} u_0 \\, , \\quad \\forall u_0 \\in H_S^\\bot \\, , \n\\end{equation}\nand \n\\begin{equation}\\label{pezzo2}\n\\Pi_S {\\cal A} u_0 = \\Pi_S ({\\cal A} - I) \\Pi_S^\\bot u_0 = \\sum_{j \\in S} \\big(u_0\\,,\\,\\psi_j \\big)_{L^2(\\mathbb T)} e^{{\\mathrm i} j x}\\,,\n\\quad \\psi_j := ({\\cal A}^T- I) e^{{\\mathrm i} j x} \\, .\n\\end{equation}\nWe claim that the difference\n\\begin{equation}\\label{differenza}\n{\\cal A} u_0 - w(1, x) = (1 + \\beta_x (x) )\\int_0^1 (1 + \\t \\b_x (x) )^{-1}\n\\big[ \\Pi_S \\partial_x (b (\\tau ) w(\\tau) ) \\big] ( \\gamma^\\t ( x + \\b (x) )) \\,d \\tau \n\\end{equation}\nwhere $ \\gamma^\\t (y) := \\gamma^{0,\\t} (y) $ is the flow of \\eqref{equazione delle caratteristiche}.\nIndeed the solution $ w(\\tau, x) $ of \\eqref{problemi di cauchy} satisfies \n$$\n\\partial_\\tau \\{ w(\\tau, \\gamma^\\t (y) ) \\} = b_{x}(\\tau, \\gamma^\\t (y)) w(\\tau, \\gamma^\\t (y)) - \\big[ \\Pi_S \\partial_x (b (\\tau ) w(\\tau) ) \\big] ( \\gamma^\\t (y))\\, .\n$$\nThen, by the variation of constant formula, we find\n$$\nw(\\tau, \\gamma^\\t (y)) = e^{\\int_0^\\tau b_x (s, \\gamma^s (y))\\,d s} \\Big( u_0(y) - \\int_0^\\tau e^{- \\int_0^s b_x(\\zeta, \n\\gamma^\\zeta (y) )\\,d\\zeta} \\big[ \\Pi_S \\partial_x (b ( s ) w(s) ) \\big] ( \\gamma^s (y)) \\, d s \\Big) \\, . \n$$\nSince $ \\partial_y \\gamma^\\t (y) $ solves the variational equation \n$ \\partial_\\tau (\\partial_y \\gamma^\\t (y)) = - b_x (\\t, \\gamma^\\t (y) ) (\\partial_y \\gamma^\\t (y)) $\nwith $ \\partial_y \\gamma^0 (y) = 1 $ we have that \n\\begin{equation}\\label{identita equazione variazionale}\ne^{\\int_0^\\tau b_x(s , \\gamma^s (y) )d s} = \\big( {\\partial_{y} \\gamma^\\t (y) } \\big)^{-1} = 1 + \\tau \\beta_x (x)\n \\end{equation}\n by remark \\ref{rem:ca}, and so we derive the expression \n$$\nw(\\tau, x) = (1 + \\tau \\beta_x (x)) \\Big\\{ u_0(x + \\tau \\beta(x)) - \\int_0^\\tau \n( 1+ s \\b_x (x) )^{-1} \n \\big[ \\Pi_S \\partial_x (b ( s ) w(s) ) \\big] ( \\gamma^s ( x + \\t \\b (x) )) \\,d s \\Big\\} \\, .\n$$\nEvaluating at $ \\t = 1 $, formula \\eqref{differenza} follows. \nNext, we develop (recall $ w(\\t) = \\Phi^\\tau (u_0 ) $) \n\\begin{equation}\\label{espressione esplicita gj}\n [\\Pi_S \\partial_x (b(\\tau) w(\\tau))] (x) \n = \\sum_{j \\in S} \\big(u_0, g_j(\\tau) \\big)_{L^2(\\mathbb T)} e^{{\\mathrm i} j x} \\, , \\quad\n g_j(\\tau) := - (\\Phi^\\tau)^T[b(\\tau) \\partial_x e^{{\\mathrm i} j x}]\\,, \n\\end{equation}\nand \\eqref{differenza} becomes \n\\begin{equation}\\label{formula utile 4}\n{\\cal A} u_0 - w(1, \\cdot) = - \\int_0^1 \\sum_{j \\in S} \\big(u_0\\,,\\, g_j(\\tau) \\big)_{L^2(\\mathbb T)} \\chi_j(\\tau, \\cdot )\\,d\\tau\\,,\n\\end{equation}\nwhere \n\\begin{equation}\\label{espressione chi j}\n\\chi_j(\\tau, x) := - ( 1 + \\b_x (x) ) (1 + \\t \\b_x (x))^{-1} e^{{\\mathrm i} j \\gamma^\\t (x + \\b(x) )} \\, . \n\\end{equation}\n\nBy \\eqref{differenza voluta}, \\eqref{pezzo2}, \\eqref{differenza}, \\eqref{formula utile 4} we deduce that \n$ \\Phi = {\\cal A}_\\bot + {\\cal R}_\\Phi $ as in \\eqref{forma buona resto cambio di variabile hamiltoniano}.\n\n\\smallskip\n\nWe now prove the estimates \\eqref{stime forma buona resto cambio di variabile hamiltoniano}. \nEach function $ \\psi_j $ in \\eqref{pezzo2} satisfies $ \\|\\psi_j \\|_s \\leq_s \\| \\beta \\|_{W^{s, \\infty}} $, \nsee \\eqref{cambio di variabile inverso}.\nThe bound $ \\| \\chi_j(\\tau)\\|_s \\leq_s 1 + \\| \\beta \\|_{W^{s + 1, \\infty}} $ follows by \\eqref{espressione chi j}.\nThe tame estimates for $ g_j(\\tau) $ defined in \\eqref{espressione esplicita gj} are more difficult because \nrequire tame estimates for the adjoint $(\\Phi^\\tau)^T$, $ \\forall \\tau \\in [0, 1] $.\nThe adjoint of the flow map can be represented \nas the flow map of the ``adjoint'' PDE \n\\begin{equation}\\label{equazione aggiunta}\n\\partial_\\tau z = \\Pi_S^\\bot \\{ b(\\tau, x) \\partial_x \\Pi_S^\\bot z \\} = \nb(\\tau, x) \\partial_x z - \\Pi_S (b(\\tau, x) \\partial_x z ) \\, , \\quad z \\in H_S^\\bot \\, , \n\\end{equation}\nwhere $ - \\Pi_S^\\bot b(\\tau,x) \\partial_x $ is the $ L^2 $-adjoint of the Hamiltonian \nvector field in \\eqref{problemi di cauchy}.\nWe denote by $ \\Psi^{\\tau_0, \\tau} $ \nthe flow of \\eqref{equazione aggiunta}, namely $ \\Psi^{\\tau_0, \\tau} (v) $ \nis the solution of \n\\eqref{equazione aggiunta} with $ \\Psi^{\\tau_0, \\tau_0} (v) = v $. \nSince the derivative $ \\partial_\\t (\\Phi^\\t (u_0) , \\Psi^{\\tau_0,\\t} (v) )_{L^2(\\mathbb T)} = 0 $, $ \\forall \\t $, we deduce that \n$ ( \\Phi^{\\tau_0} (u_0) , \\Psi^{\\tau_0,\\tau_0} (v) )_{L^2(\\mathbb T)} = ( \\Phi^0 (u_0) , \\Psi^{\\tau_0,0} (v) )_{L^2(\\mathbb T)} $, namely\n(recall that $ \\Psi^{\\tau_0,\\tau_0} (v) = v $) the adjoint \n\\begin{equation}\\label{adjoint flow}\n(\\Phi^{\\tau_0})^T = \\Psi^{\\tau_0, 0} \\,, \\quad \\forall \\tau_0 \\in [0,1] \\, . \n\\end{equation}\nThus it is sufficient to prove tame estimates for the flow $ \\Psi^{\\t_0, \\tau} $.\nWe first provide a useful expression for the solution $ z (\\tau, x) := \\Psi^{\\tau_0, \\tau} (v) $ \nof \\eqref{equazione aggiunta}, obtained by the methods of characteristics. \nLet $ \\gamma^{\\t_0,\\tau} (y) $ be the flow of \\eqref{equazione delle caratteristiche}.\nSince $ \\partial_\\tau z (\\tau, \\gamma^{\\t_0, \\t} (y)) = - [\\Pi_S (b(\\t) \\partial_x z (\\t) ] ( \\gamma^{\\t_0,\\t} (y)) $ we get \n$$\nz(\\tau, \\gamma^{\\t_0,\\tau} (y)) = v(y) + \\int_\\tau^{\\tau_0} \n[\\Pi_S (b(s) \\partial_x z (s) ] ( \\gamma^{\\t_0,s} (y)) \\,d s\\, , \\quad \\forall \\t \\in [0,1] \\, . \n$$\nDenoting by $ y = x + \\sigma(\\tau, x)$ \nthe inverse diffeomorphism of $ x = \\gamma^{\\tau_0, \\tau} (y) = y + {\\tilde \\sigma}(\\tau, y) $, we get\n\\begin{align}\n\\Psi^{\\tau_0, \\tau}(v) = z ( \\tau, x) & = \nv(x + \\sigma(\\tau, x)) + \\int_\\tau^{\\tau_0} \n[\\Pi_S (b(s) \\partial_x z (s) ] ( \\gamma^{\\t_0,s} (x + \\sigma(\\tau, x))) \\, d s \\nonumber \\\\\n& = v(x + \\sigma(\\tau, x)) +\\int_\\tau^{\\tau_0} \\sum_{j \\in S} (z(s), p_j(s)) \\kappa_j(s, x)\\,d s = v( x + \\sigma(\\tau,x) ) + {\\cal R}_\\tau v\\,, \\label{Psi-expression} \n\\end{align}\nwhere \n$ p_j (s) := - \\partial_x (b(s) e^{{\\mathrm i} j x}) $, \n$ \\kappa_j(s, x) := e^{{\\mathrm i} j \\gamma^{\\tau_0, s} (x + \\sigma(\\tau, x))} $ and \n$$\n({\\cal R}_\\tau v)(x) := \\int_\\tau^{\\tau_0} \\sum_{j \\in S} (\\Psi^{\\tau_0,s}(v), p_j(s))_{L^2(\\mathbb T)} \\kappa_j(s, x)\\,d s\\,.\n$$\nSince \n$\\| \\sigma(\\tau, \\cdot) \\|_{W^{s,\\infty}} $, $ \\| \\tilde \\sigma(\\tau, \\cdot)\\|_{W^{s,\\infty}} \n\\leq_s \\| \\beta \\|_{W^{s + 1,\\infty}}$ (recall also \\eqref{transport-free}), we derive \n$ \\| p_j \\|_s \\leq_s \\| \\beta \\|_{W^{s + 2,\\infty}} $, \n$ \\| \\kappa_j\\|_s \\leq_s 1 + \\| \\beta\\|_{W^{s + 1,\\infty}} $\nand \n$ \\| v(x + \\sigma(\\tau, x)) \\|_s \\leq_s \\| v\\|_s + \\| \\beta\\|_{W^{s + 1,\\infty}} \\| v\\|_{s_0} $, $ \\forall \\tau \\in [0,1] $. \nMoreover\n$$\n\\| {\\cal R}_\\tau v\\|_s \\leq_s {\\rm sup}_{\\tau \\in [0, 1]} \\| \\Psi^{\\tau_0,\\tau}(v) \\|_s \n\\| \\beta\\|_{W^{s_0 + 2,\\infty}} +{\\rm sup}_{\\tau \\in [0, 1]} \\| \\Psi^{\\tau_0, \\tau}(v) \\|_{s_0} \\| \\beta\\|_{W^{s + 2,\\infty}}\\, .\n$$\nTherefore, for all $ \\tau \\in [0, 1] $, \n\\begin{equation}\\label{tame aggiunto flusso parziale}\n\\| \\Psi^{\\tau_0,\\tau} v\\|_s \\leq_s \\| v\\|_s + \\| \\beta\\|_{W^{s + 1,\\infty}} \n\\| v\\|_{s_0} + {\\rm sup}_{\\tau \\in [0, 1]} \\big\\{ \\| \\Psi^{\\tau_0, \\tau} v\\|_s \\| \\beta\\|_{W^{s_0 + 2,\\infty}} +\n\\| \\Psi^{\\tau_0, \\tau} v\\|_{s_0} \\| \\beta\\|_{W^{s + 2,\\infty}} \\big\\}\\,. \n\\end{equation}\nFor $s = s_0$ it implies \n$$\n{\\rm sup}_{\\tau \\in [0, 1]} \\| \\Psi^{\\tau_0, \\tau}(v)\\|_{s_0} \\leq_{s_0} \\| v \\|_{s_0}(1 + \n\\| \\beta \\|_{W^{s_0 + 1,\\infty}}) + {\\rm sup}_{\\tau \\in [0, 1]} \\| \\Psi^{\\tau_0, \\tau}(v)\\|_{s_0} \\| \\beta\\|_{W^{s_0 + 2,\\infty}} \n$$\nand so, for $\\| \\beta\\|_{W^{s_0 + 2,\\infty} } \\leq c(s_0) $ small enough, \n\\begin{equation}\\label{tame norma bassa aggiunto flusso}\n{\\rm sup}_{\\tau \\in [0, 1]} \\| \\Psi^{\\tau_0, \\tau}(v)\\|_{s_0} \\leq_{s_0} \\| v \\|_{s_0} \\, .\n\\end{equation}\nFinally \\eqref{tame aggiunto flusso parziale}, \\eqref{tame norma bassa aggiunto flusso} imply \nthe tame estimate\n\\begin{equation}\\label{tame aggiunto flusso}\n{\\rm sup}_{\\tau \\in [0, 1]} \\| \\Psi^{\\tau_0, \\tau} (v)\\|_s \\leq_s \\| v \\|_s + \\| \\beta\\|_{W^{s + 2,\\infty}} \\| v \\|_{s_0}\\,.\n\\end{equation}\nBy \\eqref{adjoint flow} and \\eqref{tame aggiunto flusso} \nwe deduce the bound \\eqref{stime forma buona resto cambio di variabile hamiltoniano} \nfor $ g_j $ defined in \\eqref{espressione esplicita gj}.\nThe tame estimate \\eqref{stime Phi Phi -1} for $\\Phi $ follows by that of $ {\\cal A} $ and \n\\eqref{stime forma buona resto cambio di variabile hamiltoniano}\n (use Lemma \\ref{lemma:utile}). \nThe estimate for $\\Phi^{- 1}$ follows in the same way because \n$\\Phi^{-1 } = \\Phi^{1,0} $ is the backward flow.\n\\end{proof}\n\nWe conjugate $ {\\cal L}_\\omega $ in \\eqref{Lom KdVnew} \nvia the symplectic map $ \\Phi = {\\cal A}_\\bot + {\\cal R}_\\Phi $\nof Lemma \\ref{modifica simplettica cambio di variabile}. We compute \n (split $ \\Pi_S^{\\bot} = I - \\Pi_S $)\n\\begin{equation}\\label{LAbot}\n{\\cal L}_\\omega \\Phi = \\Phi {\\cal D}_\\omega + \n\\Pi_S^\\bot {\\cal A} \\big( \nb_3 \\partial_{yyy} + b_2 \\partial_{yy} + b_1 \\partial_{y} + b_0 \\big) \\Pi_S^\\bot + {\\cal R}_I \\,,\n\\end{equation}\nwhere the coefficients are\n\\begin{align}\n& b_3 (\\varphi,y) := {\\cal A}^T [ a_1 ( 1 + \\b_x)^3 ] \\label{step1: b3KdV} \n\\qquad \\qquad\nb_2 (\\varphi,y) := {\\cal A}^T \\big[ 2 (a_1)_x (1 + \\beta_x )^2 + \n6 a_1 \\beta_{xx} (1 + \\beta_x )\\big] \\\\ \n& b_1 (\\varphi,y) := \n {\\cal A}^T \\Big[ ({\\cal D}_\\omega \\beta) + 3 a_1 \\frac{\\beta_{xx}^2 }{1 + \\beta_x} + \n 4 a_1 \\beta_{xxx} + 6 (a_1)_x \\beta_{xx} + \n (a_1)_{xx} (1 + \\beta_x) + a_0 (1 + \\b_x) \\Big] \\label{tilde b1KdV} \n\\\\ \n& b_0 (\\varphi,y) := \n{\\cal A}^T \\Big[ \\frac{({\\cal D}_\\omega \\beta_x)}{1 + \\b_x} + \n a_1 \\frac{\\beta_{xxxx}}{1+ \\b_x} + \n 2 ( a_{1})_{x} \\frac{\\beta_{xxx}}{1+ \\b_x} + ( a_{1})_{xx} \\frac{\\beta_{xx}}{1+ \\b_x} +\na_0 \\frac{\\beta_{xx}}{1+ \\b_x} + (a_0)_x \\Big] \\label{tilde b0KdV}\n\\end{align}\n and the remainder\n\\begin{align}\n {\\cal R}_I &:= -\n\\Pi_S^\\bot \\partial_x ( \\varepsilon^2 {\\cal R}_2 + \\mathcal{R}_{*} ) {\\cal A}_\\bot - \\Pi_S^\\bot\n \\big( a_1 \\partial_{xxx} + 2 \n(a_1)_x \\partial_{xx} + ( (a_{1})_{xx} + a_0)\\partial_x \n+ (a_0)_x \\big) \\Pi_{S} {\\cal A} \\Pi_S^\\bot \\, \\nonumber\\\\\n & \\quad +[{\\cal D}_\\omega, {\\cal R}_\\Phi] + ({\\cal L}_\\omega - {\\cal D}_\\omega) {\\cal R}_\\Phi\\,. \\label{calR1KdV}\n\\end{align}\nThe commutator $[{\\cal D}_\\omega, {\\cal R}_\\Phi] $ has the form \\eqref{forma buona resto cambio di variabile hamiltoniano} with ${\\cal D}_\\omega g_j$ or ${\\cal D}_\\omega \\chi_j$, ${\\cal D}_\\omega \\psi_j$ instead of $\\chi_j$, $g_j$, $\\psi_j$ respectively. Also the last term $({\\cal L}_\\omega - {\\cal D}_\\omega) {\\cal R}_\\Phi$ in \\eqref{calR1KdV} has the form \\eqref{forma buona resto cambio di variabile hamiltoniano} (note that ${\\cal L}_\\omega - {\\cal D}_\\omega$ does not contain derivatives with respect to $\\varphi$). By \\eqref{LAbot}, and decomposing $ I = \\Pi_S + \\Pi_S^\\bot $, we get \n\\begin{align} \\label{L A bot finaleKdV}\n{\\cal L}_\\omega \\Phi \n= {} & \\Phi ( {\\cal D}_\\omega + b_3 \\partial_{yyy} + b_2 \\partial_{yy} + b_1 \\partial_{y} + b_0 ) \\Pi_S^\\bot\n+ {\\cal R}_ {II} \\,,\n\\\\\n\\label{calR2KdV}\n{\\cal R}_{II} \n:= {} & \\big\\{\\Pi_S^\\bot ({\\cal A} - I) \\Pi_{S} - {\\cal R}_\\Phi \\big\\} \n( b_3 \\partial_{yyy} + b_2 \\partial_{yy} + b_1 \\partial_{y} + b_0 ) \\Pi_S^\\bot + {\\cal R}_I \\,. \n\\end{align}\nNow we choose the function $ \\beta = \\b (\\varphi, x) $ such that \n\\begin{equation}\\label{choice beta}\na_1(\\varphi, x) (1 + \\b_x (\\varphi, x))^3 = b_3 (\\varphi) \n\\end{equation}\nso that the coefficient $ b_3 $ in \\eqref{step1: b3KdV} depends only on $ \\varphi $\n(note that $ {\\cal A}^T [b_3 (\\varphi)] = b_3 (\\varphi) $). \nThe only solution of \\eqref{choice beta} with zero space average is (see e.g. \\cite{BBM}-section 3.1)\n\\begin{equation}\\label{sol beta}\n\\b := \\partial_x^{-1} \\rho_0 , \\quad \n\\rho_0 := b_3 (\\varphi)^{1\/3} (a_1 (\\varphi, x))^{-1\/3} - 1, \\quad\nb_3 (\\varphi) := \\Big( \\frac{1}{2 \\pi} \\int_{\\mathbb T} (a_1 (\\varphi, x))^{-1\/3} dx \\Big)^{-3}. \n\\end{equation}\nApplying the symplectic map $ \\Phi^{-1} $ in \\eqref{L A bot finaleKdV} \nwe obtain the Hamiltonian operator (see Definition \\ref{operatore Hamiltoniano})\n\\begin{equation}\\label{cal L1 Kdv}\n{\\cal L}_1 := \\Phi^{-1} {\\cal L}_\\omega \\Phi\n= \\Pi_S^\\bot \\big( \\omega \\cdot \\partial_\\varphi + b_3(\\varphi) \\partial_{yyy} \n+ b_1 \\partial_y + b_0 \\big) \\Pi_S^\\bot + {\\mathfrak R}_1 \n\\end{equation}\nwhere $ {\\mathfrak R}_1 := \\Phi^{-1} {\\cal R}_{II} $. We used that, \nby the Hamiltonian nature of $ {\\cal L}_1 $, the coefficient \n$ b_2 = 2 (b_3)_y $ (see \\cite{BBM}-Remark 3.5) and so, by the choice \\eqref{sol beta}, we have\n$ b_2 = 2 (b_3)_y = 0 $.\nIn the next Lemma we analyse the structure of the remainder ${\\mathfrak R}_1$.\n\\begin{lemma} \\label{cal R3}\nThe operator $ {\\mathfrak R}_1 $ has the form \\eqref{forma buona con gli integrali}.\n\\end{lemma} \n\n\\begin{proof}\nThe remainders $ {\\cal R}_I $ and $ {\\cal R}_{II} $ have the form \\eqref{forma buona con gli integrali}. Indeed\n $ {\\cal R}_2, {\\cal R}_* $ in \\eqref{calR1KdV} have the form \\eqref{forma buona resto}\n (see Proposition \\ref{prop:lin}) and the term\n$ \\Pi_S {\\cal A} w = \\sum_{j \\in S } ( {\\cal A}^T e^{{\\mathrm i} j x}, w)_{L^2(\\mathbb T) } e^{{\\mathrm i} jx} $\nhas the same form. By \\eqref{forma buona resto cambio di variabile hamiltoniano}, the terms of ${\\cal R}_I$, ${\\cal R}_{II}$ which involves the operator ${\\cal R}_\\Phi$ have the form \\eqref{forma buona con gli integrali}.\nAll the operations involved preserve this structure: \nif $R_\\tau w = \\chi(\\tau) (w, g(\\tau) )_{L^2(\\mathbb T)} $, $\\tau \\in [0, 1]$, then \n\\begin{alignat*}{3}\nR_\\tau \\Pi_S^\\bot w & = \\chi(\\tau) (\\Pi_S^\\bot g(\\tau) , w)_{L^2(\\mathbb T)}\\,, \\ & \nR_\\tau {\\cal A} w & = \\chi(\\tau) ({\\cal A}^T g(\\tau) , w)_{L^2(\\mathbb T)} \\,, \\ &\n\\partial_x R_\\tau w & = \\chi_x(\\tau) (g(\\tau) , w)_{L^2(\\mathbb T)} \\,, \n\\\\\n\\Pi_S^\\bot R_\\tau w & = (\\Pi_S^\\bot \\chi(\\tau)) (g(\\tau) , w)_{L^2(\\mathbb T)} \\,, \\ & \n{\\cal A} R_\\tau w & = ({\\cal A} \\chi(\\tau)) (g(\\tau) , w)_{L^2(\\mathbb T)} \\,, \\ &\n\\Phi^{-1} R_\\tau w & = (\\Phi^{-1} \\chi(\\tau)) (g(\\tau), w)_{L^2(\\mathbb T)} \n\\end{alignat*}\n(the last equality holds because $ \\Phi^{-1} ( f (\\varphi) w ) = f (\\varphi) \\Phi^{-1} ( w ) $ for all function $f(\\varphi)$). \n Hence $ {\\mathfrak R}_1 $ has the form \\eqref{forma buona con gli integrali} \n where $ \\chi_j(\\tau) \\in H_S^\\bot $ for all $\\tau \\in [0, 1]$.\n\\end{proof}\n\nWe now put in evidence the terms of order $ \\varepsilon, \\varepsilon^2, \\ldots $, in $ b_1 $, $ b_0 $, $ \\mathfrak R_1 $, recalling that $ a_1 -1 = O(\\varepsilon^3 )$ (see \\eqref{stima a1}), $ a_0 = O( \\varepsilon) $ (see \\eqref{a1p1p2}-\\eqref{p geq 4}),\nand $ \\beta = O( \\varepsilon^3) $ (proved below in \\eqref{stima beta}).\nWe expand $ b_1$ in \\eqref{tilde b1KdV} as\n\\begin{equation}\\label{b1 Kdv}\nb_1 = - \\varepsilon p_1 -\\varepsilon^2 p_2 - q_{>2} + {\\cal D}_\\omega \\beta + 4 \\b_{xxx} + (a_1)_{xx} \n + b_{1, \\geq 4}\n\\end{equation}\nwhere $ b_{1, \\geq 4} = O(\\varepsilon^4) $ is defined by difference (the precise estimate is in Lemma \\ref{lemma:stime coeff mL1}). \n \n\\begin{remark}\\label{media beta}\nThe function \n$ {\\cal D}_\\omega \\beta $ has zero average in $ x $ by \\eqref{sol beta} as well as $ (a_1)_{xx}, \\b_{xxx} $. \n\\end{remark}\n\nSimilarly, we expand $ b_0 $ in \\eqref{tilde b0KdV} as \n\\begin{equation} \\label{b0 Kdv} \nb_0 = - \\varepsilon (p_1)_x - \\varepsilon^2 (p_2)_x - (q_{>2})_x + {\\cal D}_\\omega \\beta_x + \\b_{xxxx} + b_{0, \\geq 4} \n\\end{equation}\nwhere $ b_{0, \\geq 4} = O(\\varepsilon^4) $ is defined by difference.\n\nUsing the equalities \\eqref{calR2KdV}, \\eqref{calR1KdV} and $ \\Pi_S {\\cal A} \\Pi_S^\\bot = \\Pi_S ({\\cal A} - I) \\Pi_S^\\bot $ we get\n\\begin{equation}\\label{resto1 Kdv}\n{\\mathfrak R}_1 := \\Phi^{-1} {\\cal R}_{II} = \n- \\varepsilon^2 \\Pi_S^\\bot \\partial_x {\\cal R}_2 + {\\cal R}_{*} \n\\end{equation}\nwhere $ {\\cal R}_2 $ is defined in \\eqref{R nullo1R2} and we have renamed $\\mathcal{R}_*$ \nthe term of order $ o(\\varepsilon^2) $ in $\\mathfrak{R}_1$. \nThe remainder $ {\\cal R}_{*} $ in \\eqref{resto1 Kdv} has the form \\eqref{forma buona con gli integrali}. \n\n\\begin{lemma} \\label{lemma:stime coeff mL1}\nThere is $\\sigma = \\sigma(\\t,\\nu) > 0$ such that \n\\begin{align} \\label{stima beta}\n\\| \\b \\|_s^{\\mathrm{Lip}(\\g)} \\leq_s \\varepsilon^3 (1 + \\| {\\mathfrak{I}}_\\d \\|_{s + 1}^{\\mathrm{Lip}(\\g)} ), \\qquad \n\\| \\partial_i \\b [\\widehat \\imath ] \\|_s \n& \\leq_s \\varepsilon^3 \\big( \\| \\widehat \\imath \\|_{s + \\sigma} + \\| {\\mathfrak I}_\\d \\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma}\\big)\\,, \\\\\n\\label{stima b1 b0}\n \\| b_3 - 1 \\|_s^{\\mathrm{Lip}(\\g)} \n \\leq_s \\varepsilon^4 + \\varepsilon^{b + 2} \\| {\\mathfrak{I}}_\\d \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)}, \\qquad \\| \\partial_i b_3[\\widehat \\imath ]\\|_s & \\leq_s \\varepsilon^{b + 2} \\big( \\| \\widehat \n \\imath \\|_{s + \\sigma} + \\| {\\mathfrak{I}}_\\delta\\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma}\\big)\n\\\\\n\\label{stima Di b1 b0}\n\\| b_{1, \\geq 4} \\|_s^{\\mathrm{Lip}(\\g)} + \\| b_{0, \\geq 4} \\|_s^{\\mathrm{Lip}(\\g)} & \\leq_s \\varepsilon^4 + \\varepsilon^{b + 2} \\| {\\mathfrak{I}}_\\delta\\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} \\\\\n\\label{stima Di b geq 4}\n \\| \\partial_i b_{1, \\geq 4}[\\widehat \\imath ] \\|_s + \\| \\partial_i b_{0, \\geq 4}[\\widehat \\imath ]\\|_s & \\leq_s \\varepsilon^{b + 2}\\big( \\| \\widehat \\imath \\|_{s + \\sigma} + \\| {\\mathfrak{I}}_\\delta\\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma}\\big) . \n\\end{align}\nThe transformations $\\Phi$, $\\Phi^{-1}$ satisfy\n\\begin{align}\n\\label{stima cal A bot}\n\\|\\Phi^{\\pm 1} h \\|_s^{{\\mathrm{Lip}(\\g)}} & \\leq_s \\| h \\|_{s + 1}^{{\\mathrm{Lip}(\\g)}} + \\| {\\mathfrak I}_\\delta \\|_{s + \\sigma}^{{\\mathrm{Lip}(\\g)}} \\| h \\|_{s_0 + 1}^{{\\mathrm{Lip}(\\g)}} \\\\\n\\label{stima derivata cal A bot}\n\\| \\partial_i (\\Phi^{\\pm 1}h) [\\widehat \\imath] \\|_s & \\leq_s \n\\| h \\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma} + \n\\| h\\|_{s_0 + \\sigma} \\| \\widehat \\imath \\|_{s + \\sigma} + \n\\| {\\mathfrak I}_\\delta\\|_{s + \\sigma} \\| h\\|_{s_0 + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma}\\,.\n\\end{align}\nMoreover the remainder ${\\cal R}_{*}$ has the form \\eqref{forma buona con gli integrali},\nwhere the functions $\\chi_j(\\tau)$, $g_j(\\tau)$ satisfy the estimates \n\\eqref{stima cal R*}-\\eqref{derivate stima cal R*} uniformly in $\\tau \\in [0, 1]$.\n\\end{lemma}\n\n\n\\begin{proof}\nThe estimates \\eqref{stima beta} follow by \\eqref{sol beta}, \\eqref{stima a1}, \nand the usual interpolation and tame estimates in Lemmata \n\\ref{lemma:composition of functions, Moser}-\\ref{lemma:utile} (and Lemma \\ref{stima Aep}) and \\eqref{ansatz delta}. \nFor the estimates of $ b_3 $, by \\eqref{sol beta} and \\eqref{a1p1p2} we consider \nthe function $ r_1 $ defined in \\eqref{sigma0sigma1 def}. Recalling also\n\\eqref{finito finito} and \\eqref{T0}, the function \n$$ \nr_1 (T_\\delta ) = \\varepsilon^3 (\\partial_{u_x u_x} f_5) (v_\\d, (v_\\d)_x ) + r_{1, \\geq 4} \\, , \\quad \nr_{1, \\geq 4} := r_1 (T_\\delta ) - \\varepsilon^3 (\\partial_{u_x u_x} f_5) (v_\\d, (v_\\d)_x ) \\, . \n$$\nHypothesis (${\\mathtt S}1$) implies, as in the proof of Lemma \\ref{p3 zero average}, that the space average\n$ \\int_{\\mathbb T} (\\partial_{u_x u_x} f_5) (v_\\d, (v_\\d)_x ) dx = 0 $. Hence the bound \n\\eqref{stima b1 b0} for $ b_3 - 1 $ follows.\nFor the estimates on $\\Phi$, $\\Phi^{-1} $ we apply Lemma \\ref{modifica simplettica cambio di variabile} and the estimate \\eqref{stima beta} for $\\beta$ .\nWe estimate the remainder ${\\cal R}_*$ in \\eqref{resto1 Kdv}, using \\eqref{calR1KdV}, \\eqref{calR2KdV} and\n \\eqref{stima cal R*}-\\eqref{derivate stima cal R*}.\n\\end{proof}\n\n\n\n\\subsection{Reparametrization of time}\\label{step2}\n\nThe goal of this section is to make constant the coefficient of the highest order \nspatial derivative operator $ \\partial_{yyy} $, by a quasi-periodic reparametrization of time. \nWe consider the change of variable \n$$\n(B w)(\\varphi, y) := w(\\varphi + \\omega \\alpha(\\varphi), y), \n\\qquad \n( B^{-1} h)(\\vartheta , y ) := h(\\vartheta + \\omega \\tilde{\\alpha}(\\vartheta), y)\\,,\n$$\nwhere $ \\varphi = \\vartheta + \\omega \\tilde{\\alpha}(\\vartheta )$ \nis the inverse diffeomorphism of $ \\vartheta = \\varphi + \\omega \\alpha(\\varphi) $ in $\\mathbb T^\\nu$.\nBy conjugation, the differential operators become \n\\begin{equation} \\label{anche def rho}\nB^{-1} \\omega \\cdot \\partial_\\varphi B = \\rho(\\vartheta)\\, \\omega \\cdot \\partial_{\\vartheta} ,\n\\quad \nB^{-1} \\partial_y B = \\partial_y, \n\\quad \n\\rho := B^{-1} (1 + \\omega \\cdot \\partial_{\\varphi} \\a).\n\\end{equation}\nBy \\eqref{cal L1 Kdv}, using also that $ B $ and $ B^{-1} $ commute with $ \\Pi_S^\\bot $, we get \n\\begin{equation} \\label{secondo coniugio siti normali Kdv}\n B^{-1} {\\cal L}_1 B \n= \\Pi_S^\\bot [ \\rho \\omega \\cdot \\partial_{\\vartheta} \n + (B^{-1} b_3) \\partial_{yyy} \n + ( B^{-1} b_1 ) \\partial_y + ( B^{-1} b_0 ) ] \\Pi_S^\\bot \n + B^{-1} {\\mathfrak R}_1 B.\n\\end{equation}\nWe choose $ \\alpha $ such that \n\\begin{equation}\\label{B3solu}\n(B^{-1}b_3 )(\\vartheta ) = m_3 \\rho (\\vartheta) \\,, \\quad m_3 \\in \\mathbb R \\, , \\quad \\text{namely} \\ \\ \\ \n b_3 (\\varphi) = m_3 ( 1 + \\omega \\cdot \\partial_\\varphi \\a (\\varphi) ) \n\\end{equation}\n(recall \\eqref{anche def rho}). \nThe unique solution with zero average of \\eqref{B3solu} is\n\\begin{equation} \\label{def alpha m3}\n\\a (\\varphi) := \\frac{1}{m_3} ( \\omega \\cdot \\partial_\\varphi )^{-1} ( b_3 - m_3 ) (\\varphi) , \n\\qquad \nm_3 := \\frac{1}{(2 \\pi)^\\nu} \\int_{\\mathbb T^\\nu} b_3 (\\varphi) d \\varphi \\,. \n\\end{equation}\nHence, by \\eqref{secondo coniugio siti normali Kdv}, \n\\begin{alignat}{2} \\label{L2 Kdv}\n& B^{-1} {\\cal L}_1 B = \\rho {\\cal L}_2 \\, , \\qquad &\n& {\\cal L}_2 := \\Pi_S^\\bot ( \\omega \\cdot \\partial_{\\vartheta} + m_3 \\partial_{yyy} \n+ c_1 \\partial_y + c_0 ) \\Pi_S^\\bot + {\\mathfrak R}_2\n\\\\\n\\label{tilde c1 Kdv}\n& c_1 := \\rho^{-1} (B^{-1} b_1 ) \\,, \\qquad &\n& c_0 := \\rho^{-1} (B^{-1} b_0 ) \\, , \\qquad \n{\\mathfrak R}_2 := \\rho^{-1} B^{-1}{\\mathfrak R}_1 B \\, .\n\\end{alignat}\nThe transformed operator ${\\cal L}_2$ in \\eqref{L2 Kdv} is still Hamiltonian, since the reparametrization of time preserves the Hamiltonian structure (see Section 2.2 and Remark 3.7 in \\cite{BBM}).\n\nWe now put in evidence the terms of order $ \\varepsilon, \\varepsilon^2, \\ldots $ in $ c_1, c_0 $.\nTo this aim, we anticipate the following estimates: \n$ \\rho (\\vartheta) = 1 + O(\\varepsilon^4) $, \n$ \\alpha = O( \\varepsilon^4 \\gamma^{-1}) $, \n$ m_3 = 1 + O(\\varepsilon^4) $, $ B^{-1} - I = O( \\alpha ) $ (in low norm), \nwhich are proved in Lemma \\ref{lemma:stime coeff mL2} below. \nThen, by \\eqref{b1 Kdv}-\\eqref{b0 Kdv}, we expand the functions $ c_1, c_0 $ in \\eqref{tilde c1 Kdv} as\n\\begin{equation} \\label{tilde c1 KdV}\nc_1 = - \\varepsilon p_1 - \\varepsilon^2 p_2 - B^{-1} q_{>2} + \\varepsilon ( p_1 - B^{-1} p_1) + \\varepsilon^2 ( p_2 - B^{-1} p_2 ) + {\\cal D}_\\omega \\beta + 4 \\b_{xxx} + (a_1)_{xx} + c_{1, \\geq 4} \\, , \n\\end{equation}\n\\begin{equation} \\label{tilde c0 KdV}\nc_0 = - \\varepsilon (p_1)_x - \\varepsilon^2 (p_2)_x - (B^{-1} q_{>2})_x \n+ \\varepsilon ( p_1 - B^{-1} p_1)_x + \\varepsilon^2 ( p_2 - B^{-1} p_2 )_x + ({\\cal D}_\\omega \\beta)_x + \\b_{xxxx} + c_{0, \\geq 4}\\,, \n\\end{equation}\nwhere $ c_{1, \\geq 4}, c_{0, \\geq 4} = O( \\varepsilon^4 ) $ are defined by difference.\n\n\n\n\n\\begin{remark}\nThe functions $\\varepsilon ( p_1 - B^{-1} p_1) = O( \\varepsilon^5 \\gamma^{-1} )$ and \n$\\varepsilon^2 ( p_2 - B^{-1} p_2) = O( \\varepsilon^6 \\gamma^{-1} )$, see \\eqref{commu B p1 p2}. \nFor the reducibility scheme, the terms of order $ \\partial_x^0 $\nwith size $ O(\\varepsilon^5 \\gamma^{-1}) $ are perturbative, since\n$ \\varepsilon^5 \\gamma^{-2} \\ll 1 $.\n\\end{remark}\n\n\nThe remainder $ {\\mathfrak R}_2 $ in \\eqref{tilde c1 Kdv} has still the form \\eqref{forma buona con gli integrali} and, by \\eqref{resto1 Kdv}, \n\\begin{equation}\\label{mathfrakR2}\n{\\mathfrak R}_2 := - \\rho^{-1} B^{-1} {\\mathfrak R}_1 B = \n- \\varepsilon^2 \\Pi_S^\\bot \\partial_x {\\cal R}_2 + {\\cal R}_{*} \n\\end{equation}\nwhere $ {\\cal R}_2 $ is defined in \\eqref{R nullo1R2} and we have renamed $ {\\cal R}_{*} $ the term of order $ o( \\varepsilon^2 ) $ in $\\mathfrak{R}_2$. \n\n\n\\begin{lemma} \\label{lemma:stime coeff mL2}\nThere is $ \\sigma = \\sigma(\\nu,\\t) > 0 $ (possibly larger than $ \\sigma $ in Lemma \\ref{lemma:stime coeff mL1}) such that \n\\begin{align} \\label{stima m3}\n| m_3 - 1 |^{\\mathrm{Lip}(\\g)} \\leq C \\varepsilon^4 , \\qquad \n| \\partial_i m_3 [\\widehat \\imath ]| & \\leq C \n\\varepsilon^{b+2} \\| \\widehat \\imath \\|_{s_0 + \\sigma} \n\\\\\n\\label{stima Di b3} \n\\| \\a \\|_s^{\\mathrm{Lip}(\\g)} \\leq_s \\varepsilon^4 \\gamma^{-1} + \\varepsilon^{b + 2} \\gamma^{- 1}\\| {\\mathfrak{I}}_\\d \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)}, \\qquad \n \\| \\partial_i \\a [\\widehat \\imath] \\|_s\n& \\leq_s \\varepsilon^{b + 2} \\gamma^{-1} \\big( \\| \\widehat \\imath \\|_{s + \\sigma} + \\| {\\mathfrak I}_\\d \\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma}\\big)\\,, \n\\\\\n\\label{stima c1 c0}\n \\| \\rho -1 \\|_s^{\\mathrm{Lip}(\\g)} \n \\leq_s \\varepsilon^4 + \\varepsilon^{b + 2} \\| {\\mathfrak{I}}_\\d \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} ,\\qquad \\| \\partial_i \\rho[\\widehat \\imath ]\\|_s & \\leq_s \\varepsilon^{b + 2} \\big(\\| \\widehat \n \\imath \\|_{s + \\sigma} + \\| {\\mathfrak{I}}_\\delta\\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma} \\big) \\\\\n\\label{commu B p1 p2}\n\\| p_k - B^{-1} p_k \\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^4 \\gamma^{-1} + \\varepsilon^{b + 2} \\gamma^{- 1} \\| {\\mathfrak{I}}_\\d \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} , \\quad k = 1, 2 \n\\\\\n\\label{commu B p1 p2 der}\n\\| \\partial_i (p_k - B^{-1} p_k) [\\widehat \\imath] \\|_s \n& \\leq_s \\varepsilon^{b + 2} \\gamma^{-1} \\big( \\| \\widehat \\imath \\|_{s + \\sigma} + \\| {\\mathfrak I}_\\d \\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma}\\big)\\, \n\\\\\n\\label{stima z0 - B-1 z0}\n\\| B^{-1} q_{>2} \\|_s^{\\mathrm{Lip}(\\g)} \n& \\leq_s \\varepsilon^3 + \\varepsilon^b \\| {\\mathfrak I}_\\delta\\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} ,\n\\\\\n\\label{stima derivata z0 - B-1 z0}\n\\|\\partial_i (B^{-1} q_{>2}) [\\widehat \\imath] \\|_s \n& \\leq_s \\varepsilon^b \\big( \\| \\widehat \\imath \\|_{s + \\sigma} + \\| {\\mathfrak I}_\\d \\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma} \\big) \\,.\n\\end{align}\nThe terms \n$ c_{1, \\geq 4}, c_{0, \\geq 4} $ satisfy the bounds \\eqref{stima Di b1 b0}-\\eqref{stima Di b geq 4}. \nThe transformations $B$, $B^{-1}$ satisfy the estimates \\eqref{stima cal A bot}, \\eqref{stima derivata cal A bot}.\nThe remainder $ {\\cal R}_{*} $ has the form \\eqref{forma buona con gli integrali}, and the functions $g_j(\\tau)$, $\\chi_j(\\tau)$ satisfy \nthe estimates \\eqref{stima cal R*}-\\eqref{derivate stima cal R*} for all $\\tau \\in [0,1]$.\n\\end{lemma}\n\n\\begin{proof}\n\\eqref{stima m3} follows from \\eqref{def alpha m3},\\eqref{stima b1 b0}. \nThe estimate $ \\| \\a \\|_s \\leq_s \\varepsilon^4 \\gamma^{-1} + \\varepsilon^{b + 2} \\gamma^{- 1} \\| {\\mathfrak{I}}_\\d \\|_{s + \\sigma} $ and the inequality for \n$ \\partial_i \\alpha $ in \\eqref{stima Di b3} follow \nby \\eqref{def alpha m3},\\eqref{stima b1 b0},\\eqref{stima m3}. \nFor the first bound in \\eqref{stima Di b3} we also differentiate \n\\eqref{def alpha m3} \nwith respect to the parameter $ \\omega $. The estimates for $\\rho$ follow from $\\rho - 1 = B^{-1}(b_3 - m_3) \/ m_3$.\n\\end{proof}\n\n\\subsection{Translation of the space variable}\\label{step3}\n\nIn view of the next linear Birkhoff normal form steps (whose goal is to eliminate the terms of size $ \\varepsilon $ and $ \\varepsilon^2 $), \nin the expressions \\eqref{tilde c1 KdV}, \\eqref{tilde c0 KdV} we split \n$ p_1 = {\\bar p}_1 + ( p_1 - {\\bar p}_1)$, $p_2 = {\\bar p}_2 + ( p_2 - {\\bar p}_2)$ (see \\eqref{a1p1p2}), where\n\\begin{equation}\\label{def bar pi}\n{\\bar p}_1 := 6 {\\bar v}, \\qquad \n{\\bar p}_2 := 6 \\pi_0 [ (\\partial_x^{-1} {\\bar v})^2 ], \\qquad \n{\\bar v} (\\varphi, x) := {\\mathop \\sum}_{j \\in S} \\sqrt{\\xi_j} e^{{\\mathrm i} \\ell (j) \\cdot \\varphi } e^{{\\mathrm i} j x}, \n\\end{equation}\nand $\\ell : S \\to \\mathbb Z^\\nu$ is the odd injective map (see \\eqref{tang sites})\n\\begin{equation}\\label{del ell}\n\\ell : S \\to \\mathbb Z^\\nu, \\quad \n\\ell(\\bar \\jmath_i) := e_i, \\quad \n\\ell(-\\bar \\jmath_i) := - \\ell(\\bar \\jmath_i) = - e_i, \\quad \ni = 1,\\ldots,\\nu,\n\\end{equation}\ndenoting by $e_i = (0,\\ldots,1, \\ldots,0)$ the $i$-th vector of the canonical basis of $\\mathbb R^\\nu$.\n\n\\begin{remark}\\label{remarkp1p2}\nAll the functions $ {\\bar p}_1 $, $ {\\bar p}_2 $, $ p_1 - {\\bar p}_1 $, $ p_2 - {\\bar p}_2 $ \nhave zero average in $ x $.\n\\end{remark}\n\nWe write the variable coefficients $c_1, c_0$ of the operator $\\mathcal{L}_2$ in \\eqref{L2 Kdv} \n(see \\eqref{tilde c1 KdV}, \\eqref{tilde c0 KdV}) as\n\\begin{equation}\\label{tilde c1 KdVbis} \nc_1 = - \\varepsilon {\\bar p}_1 - \\varepsilon^2 {\\bar p}_2 + q_{c_1} + c_{1, \\geq 4}\\, , \n\\qquad \nc_0 = - \\varepsilon ({\\bar p}_1)_x - \\varepsilon^2 ({\\bar p}_2)_x + q_{c_0} + c_{0, \\geq 4}\\, , \n\\end{equation}\nwhere we define \n\\begin{gather} \\label{def qc1 qc0}\nq_{c_1} := q + 4 \\b_{xxx} + (a_1)_{xx} \\,, \\quad \nq_{c_0} := q_x + \\b_{xxxx}, \n\\\\\n\\label{def q}\nq := \\varepsilon ( p_1 - B^{-1} p_1) + \\varepsilon ( {\\bar p}_1 - p_1)\n+ \\varepsilon^2 ( p_2 - B^{-1} p_2 ) + \\varepsilon^2 ( {\\bar p}_2 - p_2)\n- B^{-1} q_{>2} + {\\cal D}_\\omega \\beta \\,.\n\\end{gather}\n\n\\begin{remark} \\label{rem:qc1 qc0}\nThe functions $ q_{c_1}, q_{c_0} $ have zero average in $ x $ (see Remarks \\ref{remarkp1p2}, \\ref{media beta} and Lemma \\ref{p3 zero average}). \n\\end{remark}\n\n\n\\begin{lemma} \\label{lemma:stime scorporo}\n The functions $\\bar p_k - p_k$, $k=1,2$ and $q_{c_m}$, \n$m = 0,1, $ satisfy\n\\begin{alignat}{2} \n\\label{stima bar pk - pk}\n& \\| {\\bar p}_k - p_k \\|_s^{\\mathrm{Lip}(\\g)} \n\\leq_s \\| {\\mathfrak{I}}_\\d \\|_{s}^{\\mathrm{Lip}(\\g)} , \n\\quad & \n& \\| \\partial_i ({\\bar p}_k - p_k)[ \\widehat \\imath] \\|_s \n\\leq_s \\| \\widehat \\imath \\|_{s} + \\| {\\mathfrak{I}}_\\d \\|_{s} \\| \\widehat \\imath \\|_{s_0}\\,, \n\\\\\n\\label{stima q}\n& \\| q_{c_m} \\|_s^{\\mathrm{Lip}(\\g)} \n\\leq_s \\varepsilon^{5}\\gamma^{-1} + \\varepsilon \\| {\\mathfrak I}_\\delta\\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)}\\,, \n\\quad & \n& \\| \\partial_i q_{c_m}[\\widehat \\imath ] \\|_s^{\\mathrm{Lip}(\\g)}\n\\leq_s \\varepsilon \\big( \\| \\widehat \\imath \\|_{s + \\sigma} + \\|{\\mathfrak I}_\\delta \\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma} \\big)\\,. \n\\end{alignat}\n\\end{lemma}\n\n\\begin{proof}\nThe bound \\eqref{stima bar pk - pk} follows from \n\\eqref{def bar pi}, \\eqref{a1p1p2}, \\eqref{T0}, \\eqref{ansatz delta}. Then use \\eqref{stima bar pk - pk},\n\\eqref{commu B p1 p2}-\\eqref{stima derivata z0 - B-1 z0},\n\\eqref{stima beta},\n\\eqref{stima a1} to prove \\eqref{stima q}. The biggest term comes from $ \\varepsilon ({\\bar p}_1 - p_1 ) $. \n\\end{proof}\n\nWe now apply the transformation $\\mathcal{T}$ defined in \\eqref{gran tau} \nwhose goal is to remove the space average from the coefficient in front of $ \\partial_y $. \n\nConsider the change of the space variable $ z = y + p(\\vartheta ) $ which induces on \n$ H^s_{S^\\bot} (\\mathbb T^{\\nu+1}) $ the operators\n\\begin{equation}\\label{gran tau}\n ({\\cal T} w)(\\vartheta, y ) := w(\\vartheta , y + p(\\vartheta)) \\, , \\quad \n ({\\cal T}^{-1} h) (\\vartheta ,z ) = h(\\vartheta, z - p(\\vartheta))\n \\end{equation}\n(which are a particular case of those used in section \\ref{step1}). \nThe differential operator becomes \n$ {\\cal T}^{-1} \\omega \\cdot \\partial_{\\vartheta} {\\cal T} $ \n$ = \\omega \\cdot \\partial_{\\vartheta} \n+ \\{ \\omega \\cdot \\partial_{\\vartheta}p (\\vartheta) \\} \\partial_z $, \n$ {\\cal T}^{-1} \\partial_{y} {\\cal T} = \\partial_{z} $. \nSince $\\mathcal{T}, \\mathcal{T}^{-1}$ commute with $ \\Pi_S^\\bot $, we get\n\\begin{align} \\label{L3 KdV}\n\\mathcal{L}_3 & := {\\cal T}^{-1}{\\cal L}_2 {\\cal T} = \\Pi_S^\\bot \\big(\\omega \\cdot \\partial_{\\vartheta} \n + m_3 \\partial_{zzz} + d_1 \\partial_z + d_0 \\big) \\Pi_S^\\bot \n + {\\mathfrak R}_3 \\,,\n\\\\\n\\label{d1d0R3}\nd_1 & := ({\\cal T}^{-1} c_1) + \\omega \\cdot \\partial_{\\vartheta} p \\, , \n\\qquad \nd_0 := {\\cal T}^{-1} c_0 \\, , \n\\qquad \n{\\mathfrak R}_3 := {\\cal T}^{-1} {\\mathfrak R}_2 {\\cal T}.\n\\end{align}\nWe choose \n\\begin{equation} \\label{def m1 p Kdv}\nm_1 := \\frac{1}{(2\\pi)^{\\nu + 1}} \\int_{\\mathbb T^{\\nu + 1}} c_1 d\\vartheta dy \\, , \\quad \np := (\\omega \\cdot \\partial_\\vartheta)^{-1} \n\\Big( m_1 - \\frac{1}{2 \\pi} \\int_{\\mathbb T} c_1 d y \\Big) \\, ,\n\\end{equation}\nso that $\\frac{1}{2 \\pi} \\int_{\\mathbb T} d_1 (\\vartheta, z) \\, dz = m_1$ for all $\\vartheta \\in \\mathbb T^\\nu$.\nNote that, by \\eqref{tilde c1 KdVbis}, \n\\begin{equation} \\label{media migliore}\n\\int_\\mathbb T c_1(\\th,y) \\, dy = \\int_\\mathbb T c_{1, \\geq 4}(\\th,y) \\, dy\\,,\n\\quad \n\\omega \\cdot \\partial_{\\vartheta} p(\\th) = m_1 - \\frac{1}{2\\pi} \\int_\\mathbb T c_{1, \\geq 4} (\\th,y)\\, dy \n\\end{equation}\nbecause $\\bar p_1, \\bar p_2, q_{c_1}$ have all zero space-average. \nAlso note that $\\mathfrak R_3$ has the form \\eqref{forma buona con gli integrali}. \nSince ${\\cal T}$ is symplectic, the operator ${\\cal L}_3$ in \\eqref{L3 KdV} is Hamiltonian.\n\\begin{remark}\nWe require Hypothesis (${\\mathtt S}1$) so that the function $ q_{>2} $ has zero space average (see Lemma \\ref{p3 zero average}). \nIf $ q_{>2} $ did not have zero average, \nthen $ p $ in \\eqref{def m1 p Kdv} would have size $O(\\varepsilon^3 \\gamma^{-1})$ \n(see \\eqref{def p3}) and, since $\\mathcal{T}^{-1} - I = O(\\varepsilon^3 \\gamma^{-1}) $, \nthe function $ \\tilde d_0 $ in \\eqref{d0 tilde KdV} would satisfy $ \\tilde d_0 = O(\\varepsilon^4 \\gamma^{-1}) $.\nTherefore it would remain a term of order $ \\partial_x^0 $ which is not \nperturbative for the reducibility scheme of section \\ref{subsec:mL0 mL5}.\n\\end{remark}\n\nWe put in evidence the terms of size $ \\varepsilon, \\varepsilon^2 $ in $ d_0 $, $ d_1 $, $ {\\mathfrak R}_3 $. \nRecalling \\eqref{d1d0R3}, \\eqref{tilde c1 KdVbis}, we split \n\\begin{equation} \\label{d1 d0 KdV}\nd_1 = - \\varepsilon {\\bar p}_1 - \\varepsilon^2 {\\bar p}_2 + \\tilde d_{1} \\, , \\quad \nd_0 = - \\varepsilon ({\\bar p}_1)_x - \\varepsilon^2 ({\\bar p}_2)_x + \\tilde d_0 \\, , \\quad \n{\\mathfrak R}_3 = \n- \\varepsilon^2 \\Pi_S^\\bot \\partial_x \\bar {\\cal R}_2 + \\widetilde \\mathcal{R}_{*}\n\\end{equation}\nwhere $ \\bar {\\cal R}_2 $ is obtained replacing \n$ v_\\d $ with ${\\bar v}$ in $ {\\cal R}_2 $ (see \\eqref{R nullo1R2}), and \n\\begin{align} \n\\label{d1 tilde KdV}\n\\tilde d_1 \n& := \\varepsilon (\\bar p_1 - {\\cal T}^{-1}{\\bar p}_1 ) \n+ \\varepsilon^2 (\\bar p_2 - {\\cal T}^{-1}{\\bar p}_2 ) \n+ {\\cal T}^{-1} ( q_{c_1} + c_{1, \\geq 4}) \n+ \\omega \\cdot \\partial_{\\vartheta} p , \n\\\\\n\\label{d0 tilde KdV}\n\\tilde d_0 \n& := \\varepsilon ( \\bar p_1 - {\\cal T}^{-1}{\\bar p}_1 )_x \n+ \\varepsilon^2 ( \\bar p_2 - {\\cal T}^{-1} {\\bar p}_2 )_x\n+ {\\cal T}^{-1} ( q_{c_0} + c_{0, \\geq 4} ),\n\\\\\n\\label{tilde mR 3}\n\\widetilde \\mathcal{R}_{*} \n& := \\mathcal{T}^{-1} \\mathcal{R}_{*} \\mathcal{T} \n+ \\varepsilon^2 \\Pi_S^\\bot \\partial_x (\\mathcal{R}_2 - \\mathcal{T}^{-1} \\mathcal{R}_2 \\mathcal{T}) + \\varepsilon^2 \\Pi_S^\\bot \\partial_x ( \\bar {\\cal R}_2 - {\\cal R}_2 ),\n\\end{align}\nand $ \\mathcal{R}_{*} $ is defined in \\eqref{mathfrakR2}. \nWe have also used that ${\\cal T}^{-1}$ commutes with $\\partial_x$ and with $\\Pi_S^\\bot$. \n\n\\begin{remark}\\label{d1 media}\nThe space average \n$ \\frac{1}{2 \\pi} \\int_{\\mathbb T} \\tilde d_1 (\\vartheta, z) \\, dz = \\frac{1}{2 \\pi} \\int_{\\mathbb T} \nd_1 (\\vartheta, z) \\, dz = m_1$ for all $\\vartheta \\in \\mathbb T^\\nu$.\n\\end{remark}\n\n\n\\begin{lemma} \\label{lemma:stime coeff mL3}\n There is $ \\sigma := \\sigma(\\nu,\\t) > 0 $ (possibly larger than in Lemma \\ref{lemma:stime coeff mL2}) such that \n\\begin{alignat}{2} \\label{stima m1}\n| m_1 |^{\\mathrm{Lip}(\\g)} & \\leq C \\varepsilon^4 , \n\\quad & | \\partial_i m_1 [\\widehat \\imath ]| & \\leq C \\varepsilon^{b+2} \\| \\widehat{\\imath} \\|_{s_0 + \\sigma} \n\\\\ \\label{stima p} \n\\| p \\|_s^{\\mathrm{Lip}(\\g)} & \\leq_s \\varepsilon^4 \\gamma^{-1} + \\varepsilon^{b + 2} \\gamma^{- 1} \\| {\\mathfrak{I}}_\\d \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)}\\,, \n\\quad & \\| \\partial_i p [\\widehat \\imath] \\|_s\n& \\leq_s \\varepsilon^{b + 2} \\gamma^{-1} \\big( \\| \\widehat \\imath \\|_{s + \\sigma} + \\| {\\mathfrak I}_\\d \\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma}\\big)\\,, \n\\\\ \\label{tilde d1 d0 KdV}\n\\| \\tilde d_k \\|_s^{\\mathrm{Lip}(\\g)} & \\leq_s \\varepsilon^{5} \\gamma^{-1} + \\varepsilon \\| {\\mathfrak{I}}_\\d \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} \\,, \n\\quad & \\| \\partial_i \\tilde d_k [\\widehat \\imath] \\|_s \n& \\leq_s \\varepsilon \\big(\\| \\widehat \\imath \\|_{s + \\sigma} + \\| {\\mathfrak I}_\\d \\|_{s + \\sigma} \n\\| \\widehat \\imath \\|_{s_0 + \\sigma} \\big)\n\\end{alignat}\nfor $k = 0,1$. Moreover the matrix $ s $-decay norm (see \\eqref{matrix decay norm}) \n\\begin{align}\\label{nuove R3}\n|\\widetilde{\\cal R}_{*}|_s^{{\\mathrm{Lip}(\\g)}} & \\leq_s \\varepsilon^{3} + \\varepsilon^{2}\\| {\\mathfrak I}_\\delta \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} \\, , \\quad \n|\\partial_i \\widetilde{\\cal R}_{*} [\\widehat \\imath ]|_s \\leq_s \\varepsilon^{2} \\| \\widehat \\imath \\|_{s + \\sigma} + \\varepsilon^{2 b - 1} \\| {\\mathfrak I}_\\delta\\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma} \\, .\n\\end{align}\nThe transformations ${\\cal T}$, ${\\cal T}^{-1}$ satisfy \\eqref{stima cal A bot}, \\eqref{stima derivata cal A bot}.\n\\end{lemma}\n\n\n\\begin{proof}\nThe estimates \\eqref{stima m1}, \\eqref{stima p} follow by \\eqref{def m1 p Kdv},\\eqref{tilde c1 KdVbis},\\eqref{media migliore}, and the bounds for $ c_{1, \\geq 4}, c_{0, \\geq 4} $ in Lemma \\ref{lemma:stime coeff mL2}. \nThe estimates \\eqref{tilde d1 d0 KdV} follow similarly by \n\\eqref{stima q},\n\\eqref{media migliore}, \\eqref{stima p}.\nThe estimates \\eqref{nuove R3} follow because\n$ \\mathcal{T}^{-1} \\mathcal{R}_{*} \\mathcal{T} $ satisfies the bounds \\eqref{stima cal R*} \nlike $ \\mathcal{R}_{*} $ does (use Lemma \\ref{remark : decay forma buona resto} and \\eqref{stima p}) \nand \n$ |\\varepsilon^2 \\Pi_S^\\bot \\partial_x ( \\bar {\\cal R}_2 - {\\cal R}_2 )|_s^{\\mathrm{Lip}(\\g)} \\leq_s \\varepsilon^2 \\| {\\mathfrak I}_\\delta \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} $. \n\\end{proof}\nIt is sufficient to estimate $ \\widetilde \\mathcal{R}_{*} $ (which has the form \\eqref{forma buona con gli integrali})\nonly in the $ s $-decay norm (see \\eqref{nuove R3}) because the next transformations will preserve it. Such norms are used \nin the reducibility scheme of section \\ref{subsec:mL0 mL5}. \n\n\\subsection{Linear Birkhoff normal form. Step 1} \\label{BNF:step1}\n\nNow we eliminate the terms of order $ \\varepsilon$ and $ \\varepsilon^2 $ of $ {\\cal L}_3 $.\nThis step is different from the \nreducibility steps that we shall perform in section \\ref{subsec:mL0 mL5}, \nbecause the diophantine constant $ \\gamma = o(\\varepsilon^2 ) $ (see \\eqref{omdio}) and so \nterms $ O( \\varepsilon), O(\\varepsilon^2) $ are not perturbative. \nThis reduction is possible thanks to the special form of the terms \n$ \\varepsilon {\\cal B}_1$, $ \\varepsilon^2 {\\cal B}_2 $ defined in \\eqref{def cal B1 B2}: \nthe harmonics of $ \\varepsilon {\\cal B}_1$, and $ \\varepsilon^2 T $ in \\eqref{mL4 T R4}, which \ncorrespond to a possible small divisor\nare naught, see Corollary \\ref{B1 zero KdV}, and Lemma \\ref{pezzo epsilon 2 A}. \nIn this section we eliminate the term $ \\varepsilon {\\cal B}_1 $. In section \\ref{BNF:step2}\nwe eliminate the terms of order $\\varepsilon^2$. \n\nNote that, since the previous transformations $ \\Phi $, $ B $, $ {\\cal T} $ are \n$ O(\\varepsilon^4 \\gamma^{-1} ) $-close to the identity, the terms of order $\\varepsilon $ and $ \\varepsilon^2 $ in $ {\\cal L}_3 $ \nare the same as in the original linearized operator.\n\n\\smallskip\n\nWe first collect all the terms of order $ \\varepsilon $ and $ \\varepsilon^2 $ in the operator $ {\\cal L}_3 $ defined in \\eqref{L3 KdV}. \nBy \\eqref{d1 d0 KdV}, \\eqref{R nullo1R2}, \\eqref{def bar pi}\nwe have, renaming $ \\vartheta = \\varphi $, $ z = x $, \n$$\n \\mathcal{L}_3 = \\Pi_S^\\bot \\big( \\omega \\cdot \\partial_\\varphi\n + m_3 \\partial_{xxx}+ \\varepsilon {\\cal B}_1 + \\varepsilon^2 {\\cal B}_2 +\n {\\tilde d}_1 \\partial_x + {\\tilde d}_0 \\big) \\Pi_S^\\bot \n + {\\widetilde {\\cal R}}_{*} \n$$\nwhere $ {\\tilde d}_1 $, $ {\\tilde d}_0 $, $ {\\widetilde {\\cal R}}_{*} $ are defined in \\eqref{d1 tilde KdV}-\\eqref{tilde mR 3} and \n(recall also \\eqref{def pi 0})\n\\begin{equation} \\label{def cal B1 B2}\n{\\cal B}_1 h := - 6 \\partial_x ( {\\bar v} h), \n\\quad\n{\\cal B}_2 h := - 6 \\partial_x \\{ {\\bar v} \n\\Pi_S [ (\\partial_x^{-1} {\\bar v} )\\,\\partial_x^{-1} h ] \n+ h \\pi_0 [ ( \\partial_x^{-1} {\\bar v} )^2 ] \\} +\n 6 \\pi_0 \\{ (\\partial_x^{-1} {\\bar v}) \\Pi_S [ {\\bar v} h ] \\}.\n\\end{equation}\nNote that ${\\cal B}_1$ and ${\\cal B}_2$ are the linear Hamiltonian vector fields of $ H_{S}^\\bot $\ngenerated, respectively, by the Hamiltonian $ z \\mapsto 3 \\int_\\mathbb T v z^2 $ in \\eqref{H3tilde}, and \nthe fourth order Birkhoff Hamiltonian $ {\\cal H}_{4,2}$ in \\eqref{mH3 mH4} at $ v = \\bar v $.\n\nWe transform $ {\\cal L}_3 $ by a symplectic operator \n$ \\Phi_1 : H_{S^\\bot}^s(\\mathbb T^{\\nu + 1}) \\to H_{S^\\bot}^s(\\mathbb T^{\\nu + 1}) $ of the form \n\\begin{equation}\\label{Phi_1}\n\\Phi_1 := {\\rm exp}(\\varepsilon A_1) = I_{H_S^\\bot} + \\varepsilon A_1 + \\varepsilon^2 \\frac{A_1^2}{2} +\n\\varepsilon^3 \\widehat A_1, \\quad \n\\widehat A_1 := {\\mathop \\sum}_{k \\geq 3} \\frac{\\varepsilon^{k - 3}}{k !} A_1^k \\, ,\n\\end{equation}\nwhere $ A_1(\\varphi) h = {\\mathop \\sum}_{j,j' \\in S^c} ( A_1)_j^{j'}(\\varphi) h_{j'} e^{{\\mathrm i} j x} $\nis a Hamiltonian vector field.\nThe map $ \\Phi_1 $ is symplectic, because it is the time-1 flow of a Hamiltonian vector field. \nTherefore\n\\begin{align}\\label{L3 new KdV}\n& {\\cal L}_3 \\Phi_1 - \\Phi_1 \\Pi_S^\\bot ( {\\cal D}_\\omega + m_3 \\partial_{xxx} ) \\Pi_S^\\bot \\\\\n& = \\Pi_S^\\bot ( \\varepsilon \\{ {\\cal D}_\\omega A_1 + m_3 [\\partial_{xxx}, A_1] + {\\cal B}_1 \\} \n+ \\varepsilon^2 \\{ {\\cal B}_1 A_1 + {\\cal B}_2 + \\frac12 m_3 [\\partial_{xxx}, A_1^2 ] + \n\\frac12 ({\\cal D}_\\omega A_1^2)\\} + {\\tilde d}_1 \\partial_x + R_3 ) \\Pi_S^\\bot\n\\nonumber\n\\end{align}\nwhere \n\\begin{align}\\label{def R3}\nR_3 & := {\\tilde d}_1 \\partial_x (\\Phi_1 - I) \\! +\\! {\\tilde d}_0 \\Phi_1\\! +\\! \\widetilde {\\cal R}_{*} \\Phi_1\\! + \\! \\varepsilon^2 {\\cal B}_2 (\\Phi_1 - I) \\! + \\! \\varepsilon^3 \\big\\{ {\\cal D}_\\omega \\widehat A_1 \\! + \\! m_3 [\\partial_{xxx}, \\widehat A_1] \\! + \\! \\frac12 {\\cal B}_1 A_1^2 \\!+ \\! \\varepsilon {\\cal B}_1 \\widehat A_1 \\big\\}\\, . \n\\end{align}\n\n\\begin{remark}\n$ R_3 $ has no longer the form \\eqref{forma buona con gli integrali}. \nHowever $ R_3 = O( \\partial_x^0 ) $ because \n$ A_1 = O( \\partial_x^{-1} ) $ (see Lemma \\ref{lemma:Dx A bounded}), and therefore $\\Phi_1 - I_{H_S^\\bot} = O(\\partial_x^{-1})$. \nMoreover the matrix decay norm of $ R_3 $ is $ o(\\varepsilon^2) $. \n\\end{remark}\n\nIn order to eliminate the order $\\varepsilon$ from \\eqref{L3 new KdV}, we choose \n\\begin{equation}\\label{cal A 1}\n( A_1)_j^{j'}(l) := \\begin{cases}\n- \\dfrac{( {\\cal B}_1)_{j}^{j'}(l)}{{\\mathrm i} (\\omega \\cdot l + m_3( j'^3 - j^3))} & \\text{if} \\ \\bar{\\omega} \\cdot l + j'^3 - j^3 \\neq 0 \\, , \\\\\n0 & \\text{otherwise},\n\\end{cases}\n\\qquad j, j' \\in S^c, \\ l \\in \\mathbb Z^\\nu.\n\\end{equation}\nThis definition is well posed. \nIndeed, by \\eqref{def cal B1 B2} and \\eqref{def bar pi}, \n\\begin{equation}\\label{def cal B1 b}\n( {\\cal B}_1)_{j}^{j'}(l):=\n\\begin{cases}\n-6 {\\mathrm i} j \\sqrt{\\xi_{j - j'}} & \\text{if} \\ j - j' \\in S\\,, \\ \\ l = \\ell(j - j') \\\\\n0 & \\text{otherwise}.\n\\end{cases}\n\\end{equation}\nIn particular $ ( {\\cal B}_1)_{j}^{j'}(l) = 0 $ unless $ | l | \\leq 1 $. Thus, for \n$ \\bar{\\omega} \\cdot l + j'^3 - j^3 \\neq 0 $, \nthe denominators in \\eqref{cal A 1} satisfy \n\\begin{align}\n|\\omega \\cdot l +m_3( j'^3 - j^3)| \n& = | m_3 (\\bar \\omega \\cdot l + j'^3 - j^3) + ( \\omega - m_3 \\bar \\omega ) \\cdot l | \n\\nonumber \\\\ & \n\\geq |m_3| | \\bar \\omega \\cdot l + j'^3 - j^3 |\n- | \\omega - m_3 \\bar \\omega | |l| \\geq 1\/2 \\, , \\quad \\forall |l| \\leq 1 \\, , \n\\label{BNFdeno}\n\\end{align}\nfor $ \\varepsilon $ small, because the non zero integer \n$ | \\bar{\\omega} \\cdot l + j'^3 - j^3| \\geq 1 $, \\eqref{stima m3}, and $ \\omega = \\bar \\omega + O(\\varepsilon^2) $. \n\n$ A_1$ defined in \\eqref{cal A 1} is a Hamiltonian vector field as $\\mathcal{B}_1$.\n\n\\begin{remark} \n\\label{rem:Ham solving homolog}\nThis is a general fact: the denominators $ \\d_{l,j,k} := {\\mathrm i} (\\omega \\cdot l + m_3( k^3 - j^3)) $ \nsatisfy $ \\overline{ \\d_{l,j,k} } = \\d_{-l,k,j} $ and \nan operator $G(\\varphi)$ is self-adjoint if and only if its matrix elements satisfy \n$ \\overline{ G_j^k(l) } = G_k^j(-l) $, see \\cite{BBM}-Remark 4.5. In a more intrinsic way, \nwe could solve \nthe homological equation of this Birkhoff step directly for the Hamiltonian function whose flow generates $ \\Phi_1 $.\n\\end{remark}\n\n\n\n\\begin{lemma} \\label{lem:cubetto} If $j,j' \\in S^c$, $j - j' \\in S$, $l = \\ell(j - j')$, then \n$ \\bar{\\omega} \\cdot l + j'^3 - j^3 = 3 j j' (j' - j) \\neq 0 $.\n\\end{lemma}\n\n\\begin{proof}\nWe have $ \\bar \\omega \\cdot l = \\bar \\omega \\cdot \\ell( j - j' ) = (j-j')^3$ because $ j - j' \\in S$ (see \\eqref{bar omega} and \\eqref{del ell}).\nNote that $ j, j' \\neq 0 $ because $ j, j' \\in S^c $, and $ j - j' \\neq 0 $ because $ j - j' \\in S $.\n\\end{proof}\n\n\\begin{corollary}\\label{B1 zero KdV}\nLet $j, j' \\in S^c$. \nIf $\\bar{\\omega} \\cdot l + j'^3 - j^3 = 0$ then $( {\\cal B}_1)_{j}^{j'}(l) = 0$. \n\\end{corollary}\n\n\\begin{proof}\nIf $({\\cal B}_1)_{j}^{j'}(l) \\neq 0$ then $j - j' \\! \\in \\! S, l = \\ell(j - j')$ \nby \\eqref{def cal B1 b}. \nHence $\\bar{\\omega} \\cdot l + j'^3 - j^3 \\! \\neq 0$ by Lemma \\ref{lem:cubetto}.\n\\end{proof}\n\nBy \\eqref{cal A 1} and the previous corollary, the term of order $\\varepsilon$ in \\eqref{L3 new KdV} is\n\\begin{equation}\\label{primo termine BNF1}\n\\Pi_S^\\bot \\big( {\\cal D}_\\omega A_1 + m_3 [\\partial_{xxx}, A_1] + {\\cal B}_1 \\big) \\Pi_S^\\bot = 0 \\, . \n\\end{equation}\nWe now estimate the transformation $ A_1 $.\n\n\\begin{lemma}\\label{lem: A decay} \n$(i)$ For all $l \\in \\mathbb Z^\\nu$, $j,j' \\in S^c$, \n\\begin{equation}\\label{Acoef}\n| (A_1)_{j}^{j'}(l)| \\leq C (| j | + | j' |)^{-1}\\, , \\quad \n| (A_1)_{j}^{j'}(l)|^{\\rm lip} \\leq \\varepsilon^{-2} (|j| + |j'|)^{-1} \\,.\n\\end{equation}\n$(ii)$ $ (A_1)_j^{j'}(l) = 0$ for all $l \\in \\mathbb Z^\\nu$, $j,j' \\in S^c$ such that $|j - j'| > C_S $, where $C_S := \\max\\{ |j| : j \\in S \\}$. \n\\end{lemma}\n \n\n\\begin{proof}\n$(i)$ We already noted that $ (A_1)_j^{j'}(l) = 0$, $ \\forall |l| > 1$. \nSince $ | \\omega | \\leq |\\bar \\omega | + 1 $, one has, for $ |l| \\leq 1 $, $ j \\neq j' $,\n\\[\n|\\omega \\cdot l + m_3 (j'^3 - j^3)| \n\\geq |m_3||j'^3 - j^3| - |\\omega \\cdot l| \n\\geq \\frac14 (j'^2 + j^2) - |\\omega| \n\\geq \\frac18 (j'^2 + j^2) \\, , \\quad \\forall (j'^2 + j^2) \\geq C, \n\\]\nfor some constant $C > 0$. \nMoreover, recalling that also \\eqref{BNFdeno} holds, we deduce that for $ j \\neq j' $, \n\\begin{equation} \\label{lower bound}\n (A_1)_j^{j'}(l) \\neq 0 \n\\quad \\Rightarrow \\quad\n|\\omega \\cdot l + m_3 (j'^3 - j^3)| \\geq c ( | j | + | j' | )^2 \\, . \n\\end{equation}\nOn the other hand, if $ j = j' $, $ j \\in S^c$, the matrix $ (A_1)_j^j (l) = 0 $, $ \\forall l \\in \\mathbb Z^\\nu$, \nbecause $ ({\\cal B}_1)_j^j (l) = 0 $ by \\eqref{def cal B1 b} \n(recall that $ 0 \\notin S $). Hence \\eqref{lower bound} holds for all $ j, j' $.\nBy \\eqref{cal A 1}, \\eqref{lower bound}, \\eqref{def cal B1 b} we deduce the first \nbound in \\eqref{Acoef}. \nThe Lipschitz bound follows similarly (use also $ |j - j'| \\leq C_S $). \n$(ii)$ follows by \\eqref{cal A 1}-\\eqref{def cal B1 b}. \n\\end{proof}\n\nThe previous lemma means that $ A = O(| \\partial_x|^{-1})$. \nMore precisely we deduce that\n\n\\begin{lemma} \\label{lemma:Dx A bounded}\n$ | A_1 \\partial_x |_s^{\\mathrm{Lip}(\\g)} + | \\partial_x A_1 |_s^{\\mathrm{Lip}(\\g)} \\leq C(s) $.\n\\end{lemma}\n\n\\begin{proof}\nRecalling the definition of the (space-time) matrix norm in \\eqref{decayTop}, \nsince $(A_1)_{j_1}^{j_2}(l) = 0$ outside the set of indices $|l| \\leq 1, |j_1 - j_2| \\leq C_S$, \nwe have\n\\begin{align*}\n| \\partial_x A_1 |_s^2 \n& = \\sum_{|l| \\leq 1, \\, |j| \\leq C_S} \n\\Big( \\sup_{j_1 - j_2 = j} |j_1| | (A_1)_{j_1}^{j_2}(l)| \\Big)^2 \n\\langle l,j \\rangle^{2s} \n\\leq C(s) \n\\end{align*}\nby Lemma \\ref{lem: A decay}. The estimates for $ |A_1 \\partial_x |_s $ and the Lipschitz bounds \nfollow similarly. \n\\end{proof}\n\n\nIt follows that the symplectic map $ \\Phi_1 $ in \\eqref{Phi_1} is invertible for $ \\varepsilon $ small, \nwith inverse\n\\begin{equation}\\label{A1 check}\n \\Phi_1^{-1} = {\\rm exp}(-\\varepsilon A_1) = I_{H_S^\\bot} + \\varepsilon {\\check A}_1 \\, , \\ \n {\\check A}_1 := {\\mathop \\sum}_{n \\geq 1} \\frac{\\varepsilon^{n-1}}{n !} (-A_1)^n \\, , \\ \n| {\\check A}_1 \\partial_x |_s^{\\mathrm{Lip}(\\g)} + | \\partial_x {\\check A}_1 |_s^{\\mathrm{Lip}(\\g)} \\leq C(s) \\, . \n\\end{equation}\nSince $ A_1 $ solves the homological equation \\eqref{primo termine BNF1}, the $ \\varepsilon $-term in\n\\eqref{L3 new KdV} is zero, and, \nwith a straightforward calculation, \nthe $ \\varepsilon^2 $-term simplifies to $ {\\cal B}_2 + \\frac12 [{\\cal B}_1, A_1] $. \nWe obtain the Hamiltonian operator\n\\begin{align} \\label{bernardino1}\n{\\cal L}_4 & := \\Phi_1^{-1} {\\cal L}_3 \\Phi_1 \n= \\Pi_S^\\bot ( {\\cal D}_\\omega + m_3 \\partial_{xxx} + {\\tilde d}_1 \\partial_x \n+ \\varepsilon^2 \\{ {\\cal B}_2 + \\tfrac12 [{\\cal B}_1, A_1] \\} + \\tilde{R}_4 ) \\Pi_S^\\bot \n\\\\\n\\label{bernardino 2}\n{\\tilde R}_4 & := (\\Phi_1^{-1} - I) \\Pi_S^\\bot [ \\varepsilon^2 ( {\\cal B}_2 + \\tfrac12 [{\\cal B}_1, A_1] ) + {\\tilde d}_1 \\partial_x ] + \\Phi_1^{-1} \\Pi_S^\\bot R_3 \\, .\n\\end{align}\nWe split $ A_1 $ defined in \\eqref{cal A 1}, \\eqref{def cal B1 b} into $ A_1 = \\bar{A}_1 + \\widetilde{A}_1$ where, for all $ j, j' \\in S^c $, $l \\in \\mathbb Z^\\nu$, \n\\begin{equation}\\label{bar A1} \n( {\\bar A}_1)_j^{j'}(l) := \n\\dfrac{6 j \\sqrt{\\xi_{j - j'}}}{ \\bar\\omega \\cdot l + j'^3 - j^3} \n\\qquad \\text{if } \\ \\bar{\\omega} \\cdot l + j'^3 - j^3 \\neq 0, \\ \\ j-j' \\in S , \\ \\ \nl = \\ell(j - j'), \n\\end{equation}\nand $( {\\bar A}_1)_j^{j'}(l) := 0$ otherwise. \nBy Lemma \\ref{lem:cubetto}, for all $ j , j' \\in S^c $, $l \\in \\mathbb Z^\\nu$, \n$ ( {\\bar A}_1)_j^{j'}(l) = \n\\frac{2 \\sqrt{\\xi_{j - j'}}}{ j' ( j' - j )} \n$ if $ j-j' \\in S $, $ l = \\ell(j - j') $, and \n$ ( {\\bar A}_1)_j^{j'}(l) = 0 $ otherwise, \nnamely (recall the definition of $ \\bar v $ in \\eqref{def bar pi}) \n\\begin{equation}\\label{bar A1 esplicita}\n{\\bar A}_1h = 2 \\Pi_S^\\bot [( \\partial_x^{-1} {\\bar v} )(\\partial_x^{-1} h)] \\, , \n\\quad \\forall h \\in H_{S^\\bot}^s(\\mathbb T^{\\nu+1}) \\,.\n\\end{equation}\nThe difference is\n\\begin{equation}\\label{widetilde A1}\n(\\widetilde{A}_1)_j^{j'}(l) = (A_1 - \\bar{A}_1)_j^{j'}(l) = - \n\\frac{6 j \\sqrt{\\xi_{j - j'}}\\big\\{ ( \\omega - \\bar \\omega) \\cdot l + (m_3 - 1)(j'^3 - j^3) \\big\\}}{\\big( \\omega \\cdot l + m_3(j'^3 - j^3) \\big) \\big(\\bar\\omega \\cdot l + j'^3 - j^3 \\big)} \n\\end{equation}\nfor $ j, j' \\in S^c $, $ j - j' \\in S $, $ l = \\ell(j - j') $, \nand $ (\\widetilde{A}_1)_j^{j'}(l) = 0 $ otherwise. \nThen, by \\eqref{bernardino1}, \n\\begin{equation}\\label{mL4 T R4}\n{\\cal L}_4 = \\Pi_S^\\bot \\big( {\\cal D}_\\omega + m_3 \\partial_{xxx} \n+ {\\tilde d}_1 \\partial_x \n+ \\varepsilon^2 T + R_4 \\big) \\Pi_S^\\bot \\,, \n\\end{equation}\nwhere\n\\begin{equation}\\label{T hamiltoniano}\nT := {\\cal B}_2 + \\frac12 [{\\cal B}_1, {\\bar A}_1] \\,, \\qquad R_4 := \\frac{\\varepsilon^2}{2} [{\\cal B}_1, \\widetilde A_1] + \\tilde R_4\\,.\n\\end{equation}\nThe operator $ T $ is Hamiltonian as $ {\\cal B}_2 $, $ {\\cal B}_1 $, $ {\\bar A}_1 $ (the commutator of two Hamiltonian vector fields is Hamiltonian). \n\n\\begin{lemma} \\label{lemma:R4}\n There is $ \\sigma = \\sigma(\\nu,\\t) > 0 $ (possibly larger than in Lemma \\ref{lemma:stime coeff mL3}) such that \n\\begin{align} \\label{stima Lip R4} \n| R_4 |_s^{{\\rm Lip}(\\gamma)} \n& \\leq_s \\varepsilon^{5} \\gamma^{-1} + \\varepsilon \\| {\\mathfrak I}_\\d \\|_{s + \\sigma}^{{\\rm Lip}(\\gamma)} \\,, \\quad \n| \\partial_i R_4 [\\widehat \\imath] |_s \n\\leq_s \\varepsilon \\big( \\| \\widehat \\imath \\|_{s + \\sigma} + \\| {\\mathfrak I}_\\d \\|_{s + \\sigma} \n \\| \\widehat \\imath \\|_{ s_0 + \\sigma} \\big) \\, . \n \\end{align}\n\\end{lemma}\n\n\\begin{proof} \nWe first estimate $ [{\\cal B}_1, \\widetilde A_1] = \n({\\cal B}_1 \\partial_x^{-1}) (\\partial_x {\\widetilde A}_1)- ({\\widetilde A}_1 \\partial_x)( \\partial_x^{-1} {\\cal B}_1) $.\nBy \\eqref{widetilde A1}, $ | \\omega - \\bar \\omega | \\leq C \\varepsilon^2 $ \n(as $ \\omega \\in \\Omega_\\varepsilon $ in \\eqref{Omega epsilon}) and \\eqref{stima m3}, \narguing as in Lemmata \\ref{lem: A decay}, \\ref{lemma:Dx A bounded}, we deduce that \n$ | {\\widetilde A}_1 \\partial_x |_s^{\\mathrm{Lip}(\\g)} + $ $ | \\partial_x {\\widetilde A}_1 |_s^{\\mathrm{Lip}(\\g)} \\leq_s \\varepsilon^2 $.\nBy \\eqref{def cal B1 B2} \nthe norm $ |{\\cal B}_1 \\partial_x^{-1}|_s^{\\mathrm{Lip}(\\g)} + |\\partial_x^{-1} {\\cal B}_1|^{\\mathrm{Lip}(\\g)} \\leq C(s) $. \nHence $ \\varepsilon^2 |[{\\cal B}_1, {\\widetilde A}_1]|_s^{\\mathrm{Lip}(\\g)} \\leq_s \\varepsilon^4 $. \nFinally \\eqref{T hamiltoniano},\n\\eqref{bernardino 2}, \\eqref{A1 check}, \\eqref{def R3}, \\eqref{tilde d1 d0 KdV}, \\eqref{nuove R3}, \n and the interpolation estimate \\eqref{interpm Lip} imply \\eqref{stima Lip R4}.\n\\end{proof}\n\n\n\n\n\\subsection{Linear Birkhoff normal form. Step 2}\\label{BNF:step2} \n\n\nThe goal of this section is to remove the term $ \\varepsilon^2 T$\nfrom the operator $\\mathcal{L}_4$ defined in \\eqref{mL4 T R4}.\nWe conjugate the Hamiltonian operator $\\mathcal{L}_4$ via a symplectic map\n\\begin{equation}\\label{def Phi2}\n\\Phi_2 := {\\rm exp}(\\varepsilon^2 A_2) = I_{H_S^\\bot} + \\varepsilon^2 A_2 + \\varepsilon^4 \\widehat A_2\\,,\\quad \n\\widehat A_2 := {\\mathop \\sum}_{k \\geq 2} \\frac{\\varepsilon^{2(k - 2)}}{k !} A_2^k\n\\end{equation}\nwhere $A_2(\\varphi) = {\\mathop \\sum}_{j,j' \\in S^c} ( A_2)_j^{j'}(\\varphi) h_{j'} e^{{\\mathrm i} j x} $ \nis a Hamiltonian vector field. We compute \n\\begin{gather} \\label{L4 diff}\n{\\cal L}_4 \\Phi_2 - \\Phi_2 \\Pi_S^\\bot \\big( {\\cal D}_\\omega + m_3 \\partial_{xxx} \\big) \\Pi_S^\\bot \n= \\, \\Pi_S^\\bot ( \\varepsilon^2 \\{ {\\cal D}_\\omega A_2 + m_3 [\\partial_{xxx}, A_2 ] + T \\} \n+ \\tilde{d}_1 \\partial_x + \\tilde R_5 ) \\Pi_S^\\bot \\,,\n\\\\\n\\label{def tilde R5}\n\\tilde R_5 := \\Pi_S^\\bot \\{ \n\\varepsilon^4 ( ({\\cal D}_\\omega \\widehat A_2) + m_3 [\\partial_{xxx}, \\widehat A_2] )\n+ (\\tilde{d}_1 \\partial_x + \\varepsilon^2 T)(\\Phi_2 - I) + R_4 \\Phi_2 \\} \\Pi_S^\\bot \\, . \n\\end{gather}\nWe define\n\\begin{equation}\\label{A2}\n( A_2)_j^{j'}(l) := \n- \\dfrac{T_j^{j'}(l)}{ {\\mathrm i} (\\omega \\cdot l + m_3 (j'^3 - j^3) )} \n\\quad \\text{if } \\ \\bar\\omega \\cdot l + j'^3 - j^3 \\neq 0; \n\\qquad \n( A_2)_j^{j'}(l) := 0 \\quad \\text{otherwise.}\n\\end{equation}\nThis definition is well posed. \nIndeed, by \\eqref{T hamiltoniano}, \\eqref{def cal B1 b}, \\eqref{bar A1}, \\eqref{def cal B1 B2}, \nthe matrix entries $ T_j^{j'} (l) = 0 $ for all $ | j - j' | > 2 C_S $, $l \\in \\mathbb Z^\\nu $, \nwhere $C_S := \\max \\{ |j| \\, , j \\in S \\} $. \nAlso $ T_j^{j'} (l) = 0 $ for all $ j, j' \\in S^c $, $ | l | > 2 $ (see also \n\\eqref{forma funzionale B1 A1}, \\eqref{B1B2 Fourier}, \\eqref{B3 Fourier} below).\nThus, arguing as in \\eqref{BNFdeno}, if $ \\bar\\omega \\cdot l + j'^3 - j^3 \\neq 0 $, then \n$ |\\omega \\cdot l + m_3 (j'^3 - j^3)| \\geq 1 \/ 2 $.\nThe operator $ A_2 $ is a Hamiltonian vector field because $ T $ is Hamiltonian and by Remark \\ref{rem:Ham solving homolog}.\n\nNow we prove that the Birkhoff map \n$\\Phi_2$ removes completely the term $\\varepsilon^2 T$. \n\n\\begin{lemma} \\label{pezzo epsilon 2 A}\n Let $j, j' \\in S^c$. If $\\bar \\omega \\cdot l + j'^3 - j^3 = 0$, then $T_j^{j'}(l) = 0$.\n\\end{lemma}\n\n\\begin{proof}\nBy \\eqref{def cal B1 B2}, \\eqref{bar A1 esplicita} we get\n$ {\\cal B}_1 {\\bar A}_1 h \n= - 12 \\partial_x \\{ \\bar v \\Pi_S^\\bot[( \\partial_x^{-1} \\bar v )(\\partial_x^{-1} h)] \\}\n$, ${\\bar A}_1{\\cal B}_1 h = $ $ - 12 \\Pi_S^\\bot [ (\\partial_x^{-1} \\bar v) \\Pi_S^\\bot (\\bar v h) ] $ for all $h \\in H_{S^\\bot}^s $, \nwhence, recalling \\eqref{def bar pi}, for all $ j, j' \\in S^c $, $ l \\in \\mathbb Z^\\nu $, \n\\begin{equation}\\label{forma funzionale B1 A1}\n( [{\\cal B}_1, {\\bar A}_1])_{j}^{j'}(l) = 12 {\\mathrm i} \\!\\!\\! \\sum_{\\begin{subarray}{c}\nj_1, j_2 \\in S, \\, j_1 + j_2 = j - j' \\\\\nj' + j_2 \\in S^c, \\, \\ell(j_1) + \\ell(j_2) = l \n\\end{subarray}} \\!\\!\\! \\frac{ j j_1 -j' j_2 }{j' j_1 j_2} \\sqrt{\\xi_{j_1} \\xi_{j_2}}\\,, \n\\end{equation}\nIf $([{\\cal B}_1, {\\bar A}_1])_j^{j'}(l) \\neq 0 $\nthere are $ j_1, j_2 \\in S $ such that $ j_1 + j_2 = j - j' $, $j' + j_2 \\in S^c$, $ \\ell(j_1) + \\ell(j_2) = l $.\nThen\n\\begin{equation} \\label{con lo zero}\n\\bar \\omega \\cdot l + j'^3 - j^3 \n = \\bar \\omega \\cdot \\ell(j_1 ) + \\bar \\omega \\cdot \\ell(j_2) + j'^3 - j^3 \n \\stackrel{\\eqref{del ell}} \n = j_1^3 + j_2^3 + j'^3 - j^3\\, .\n\\end{equation}\nThus, if $\\bar \\omega \\cdot l + j'^3 - j^3 = 0$, \nLemma \\ref{lemma:interi} implies $ (j_1 + j_2 )( j_1 + j') ( j_2 + j' ) = 0 $. Now $ j_1 + j' $, $ j_2 + j' \\neq 0 $\nbecause $ j_1, j_2 \\in S $, $ j' \\in S^c $ and $ S $ is symmetric. \nHence $ j_1 + j_2 = 0 $, which implies $ j = j' $ and $ l = 0 $ (the map $ \\ell $ in \\eqref{del ell} is odd).\nIn conclusion, if $\\bar \\omega \\cdot l + j'^3 - j^3 = 0$, the only nonzero matrix entry \n$ ([{\\cal B}_1, {\\bar A}_1])_j^{j'}(l) $ is\n\\begin{equation}\\label{parte diagonale zero}\n([{\\cal B}_1, \\bar A_1])_j^j(0) \n\\stackrel{\\eqref{forma funzionale B1 A1}} = 24 {\\mathrm i} \\sum_{j_2 \\in S, \\, j_2 + j \\in S^c } \n\\xi_{j_2} {j_2^{-1}}.\n\\end{equation}\nNow we consider $ {\\cal B}_2 $ in \\eqref{def cal B1 B2}. Split $ {\\cal B}_2 = B_1 + B_2 + B_3 $, \nwhere \n$B_1 h := - 6 \\partial_x \\{ {\\bar v} \\Pi_S [ (\\partial_x^{-1} {\\bar v}) \\partial_{x}^{-1} h ] \\}$, \n$B_2 h := - 6 \\partial_x \\{ h \\pi_0 [(\\partial_x^{-1} {\\bar v} )^2] \\}$, \n$B_3 h := 6 \\pi_0 \\{ \\Pi_S ({\\bar v} h) \\partial_x^{-1} {\\bar v} \\}$.\nTheir Fourier matrix representation is\n\\begin{gather}\\label{B1B2 Fourier}\n(B_1)_j^{j'}(l) = 6 {\\mathrm i} j \\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\! \\sum_{\\begin{subarray}{c}\nj_1, j_2 \\in S, \\, j_1 + j' \\in S \\\\\nj_1 + j_2 = j - j', \\, \\ell(j_1) + \\ell(j_2) = l \n\\end{subarray}} \\!\\!\\!\\! \\!\\!\\!\\!\\!\\!\\! \\frac{\\sqrt{\\xi_{j_1} \\xi_{j_2}}}{ j_1 j'} \\, , \n\\qquad \n(B_2)_{j}^{j'}(l) = 6 {\\mathrm i} j \\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\! \\sum_{\\begin{subarray}{c}\nj_1 , j_2 \\in S, \\, j_1 + j_2 \\neq 0 \\\\\nj_1 + j_2 = j - j', \\, \\ell(j_1) + \\ell(j_2) = l\n\\end{subarray}} \\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\! \\frac{\\sqrt{\\xi_{j_1} \\xi_{j_2}}}{j_1 j_2} \\,, \n\\\\\n\\label{B3 Fourier}\n(B_3)_j^{j'}(l) = 6 \\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\! \\sum_{\\begin{subarray}{c}\nj_1 , j_2 \\in S, \\, j_1 + j' \\in S \\\\\nj_1 + j_2 = j - j', \\, \\ell(j_1) + \\ell(j_2) = l\n\\end{subarray}} \\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\! \\frac{\\sqrt{\\xi_{j_1} \\xi_{j_2}}}{{\\mathrm i} j_2}\\,,\n\\qquad j,j' \\in S^c, \\ l \\in \\mathbb Z^\\nu.\n\\end{gather}\nWe study the terms $ B_1 $, $ B_2 $, $ B_3$ separately.\nIf $(B_1)_j^{j'}(l) \\neq 0$, there are $j_1 , j_2 \\in S$ such that \n$ j_1 + j_2 = j - j' $, $j_1 + j' \\in S$, $ l = \\ell(j_1) + \\ell(j_2) $ and\n\\eqref{con lo zero} holds. Thus, if $\\bar \\omega \\cdot l + j'^3 - j^3 = 0$, \nLemma \\ref{lemma:interi} implies \n$ (j_1 + j_2) (j_1 + j')(j_2 + j') = 0 $, \nand, since $j' \\in S^c $ and $ S $ is symmetric, the only possibility is $j_1 + j_2 = 0$. Hence $j = j'$, $l = 0$.\nIn conclusion, if $\\bar \\omega \\cdot l + j'^3 - j^3 = 0$, the only nonzero matrix element \n$ ( B_1 )_j^{j'}(l) $ is \n\\begin{equation} \\label{diag R0}\n(B_1)_j^j(0) = 6 {\\mathrm i} \\sum_{ j_1 \\in S, \\, j_1 + j \\in S } \\xi_{j_1} j_1^{-1} \\,.\n\\end{equation}\nBy the same arguments, if $ (B_2)_j^{j'} (l) \\neq 0 $ and $\\bar \\omega \\cdot l + j'^3 - j^3 = 0$ we find $ (j_1 + j_2) (j_1 + j')(j_2 + j') = 0 $, \nwhich is impossible because also $ j_1 + j_2 \\neq 0$.\nFinally, arguing as for $ B_1 $, if $\\bar \\omega \\cdot l + j'^3 - j^3 = 0$, then the only nonzero matrix element $ ( B_3 )_j^{j'}(l) $ is\n\\begin{equation} \\label{diag R1 R2}\n(B_3)_j^j(0) = 6 {\\mathrm i} \\sum_{j_1 \\in S, \\, j_1 + j \\in S } \\xi_{j_1} j_1^{-1} \\,.\n\\end{equation}\nFrom \\eqref{parte diagonale zero}, \\eqref{diag R0}, \\eqref{diag R1 R2} we deduce that, if \n$ \\bar \\omega \\cdot l + j'^3 - j^3 = 0 $, then the only non zero elements $ (\\frac12 [\\mathcal{B}_1, \\bar A_1] + B_1 + B_3)_j^{j'} (l) $ must be for $ (l, j, j') = (0, j, j) $. \nIn this case, we get\n\\begin{equation}\\label{Tjj0}\n\\frac12( [\\mathcal{B}_1, \\bar A_1])_j^j(0) + (B_1)_j^j(0) + (B_3)_j^j(0) \n= 12 {\\mathrm i} \\sum_{\\begin{subarray}{c} j_1 \\in S \\\\ j_1 + j \\in S^c \\end{subarray}} \\frac{\\xi_{j_1}}{ j_1}\n+ 12 {\\mathrm i} \\sum_{\\begin{subarray}{c} j_1 \\in S \\\\ j_1 + j \\in S \\end{subarray}} \\frac{\\xi_{j_1}}{ j_1}\n= 12 {\\mathrm i} \\sum_{j_1 \\in S} \\frac{\\xi_{j_1}}{ j_1} \n= 0\n\\end{equation}\nbecause \nthe case $j_1 + j = 0$ is impossible ($j_1 \\in S$, $j' \\in S^c$ and $S$ is symmetric), \nand the function $S \\ni j_1 \\to \\xi_{j_1} \/ j_1 \\in \\mathbb R$ is odd. \nThe lemma follows by \\eqref{T hamiltoniano}, \\eqref{Tjj0}.\n\\end{proof}\n\nThe choice of $ A_2 $ in \\eqref{A2} and Lemma \\ref{pezzo epsilon 2 A} imply that \n\\begin{equation}\\label{pezzo zero}\n \\Pi_S^\\bot \\big( {\\cal D}_\\omega A_2 + m_3 [\\partial_{xxx}, A_2] + T \\big) \\Pi_S^\\bot = 0 \\, .\n\\end{equation}\n\n\\begin{lemma}\\label{A2 decay}\n$ |\\partial_x A_2|_s^{\\mathrm{Lip}(\\g)} + | A_2 \\partial_x|_s^{\\mathrm{Lip}(\\g)} \\leq C(s) $. \n\\end{lemma}\n\n\\begin{proof}\nFirst we prove that the diagonal elements $ T_j^j (l) = 0 $ for all $ l \\in \\mathbb Z^\\nu $. \nFor $l = 0$, we have already proved that $T_j^j(0) = 0$ (apply Lemma \\ref{pezzo epsilon 2 A} with $j = j'$, $l=0$).\nMoreover, in each term $[\\mathcal{B}_1, \\bar A_1]$, $B_1$, $B_2$, $B_3$ (see \\eqref{forma funzionale B1 A1}, \\eqref{B1B2 Fourier}, \\eqref{B3 Fourier}) \nthe sum is over $j_1 + j_2 = j - j'$, $l = \\ell(j_1) + \\ell(j_2)$. \nIf $j = j'$, then $j_1 + j_2 = 0$, and $l = 0$. \nThus $ T_j^j (l) = T_j^j (0) = 0 $. \nFor the off-diagonal terms $ j \\neq j' $ we argue as in Lemmata \\ref{lem: A decay}, \\ref{lemma:Dx A bounded},\nusing that all the denominators $ |\\omega \\cdot l + m_3 ( j'^3 - j^3 ) | \\geq c ( |j| + |j'|)^2 $. \n\\end{proof}\n\nFor $ \\varepsilon $ small, the map $ \\Phi_2 $ in \\eqref{def Phi2} is invertible and $ \\Phi_2 = \\exp(-\\varepsilon^2 A_2 ) $. \nTherefore \\eqref{L4 diff}, \\eqref{pezzo zero} imply \n\\begin{align} \\label{L5 KdV}\n{\\cal L}_5 \n& := \\Phi_2^{-1} {\\cal L}_4 \\Phi_2 \n= \\Pi_S^\\bot ( {\\cal D}_\\omega + m_3 \\partial_{xxx} + \\tilde{d}_1 \\partial_x + R_5 ) \\Pi_S^\\bot \\,,\n\\\\\n\\label{R5}\nR_5 \n& := ( \\Phi_2^{-1} - I) \\Pi_S^\\bot \\tilde{d}_1 \\partial_x \n+ \\Phi_2^{-1} \\Pi_S^\\bot {\\tilde R}_5 \\,. \n\\end{align}\nSince $ A_2 $ is a Hamiltonian vector field, the map $\\Phi_2$ is symplectic and so ${\\cal L}_5$ is Hamiltonian.\n\\begin{lemma} \\label{lemma:R5}\n$R_5$ satisfies the same estimates \\eqref{stima Lip R4} as $R_4$ (with a possibly larger $\\sigma$).\n\\end{lemma}\n\n\\begin{proof} \nUse \\eqref{R5}, Lemma \\ref{A2 decay}, \\eqref{tilde d1 d0 KdV}, \\eqref{def tilde R5}, \n\\eqref{stima Lip R4} and the interpolation inequalities \\eqref{multiplication Lip}, \\eqref{interpm Lip}. \n\\end{proof} \n\n\n\\subsection{Descent method}\\label{step5}\n\nThe goal of this section is to transform $ {\\cal L}_5 $ in \\eqref{L5 KdV} so that the coefficient of $ \\partial_x $ becomes constant.\nWe conjugate $ {\\cal L}_5 $ via a symplectic map of the form \n\\begin{equation}\\label{def descent}\n{\\cal S} := {\\rm exp}(\\Pi_S^\\bot (w \\partial_x^{-1}))\\Pi_S^\\bot = \\Pi_S^\\bot \\big( I + w \\partial_x^{-1} \\big) \\Pi_S^\\bot + \\widehat {\\cal S}\\,,\\quad \\widehat {\\cal S} := {\\mathop \\sum}_{k \\geq 2} \n\\frac{1}{k!} [\\Pi_S^\\bot (w \\partial_x^{-1})]^k \\Pi_S^\\bot\\,,\n\\end{equation}\nwhere $ w : \\mathbb T^{\\nu+1} \\to \\mathbb R $ is a function. \nNote that $\\Pi_S^\\bot (w \\partial_x^{-1}) \\Pi_S^\\bot$ is the Hamiltonian vector field generated by \n $ - \\frac12 \\int_\\mathbb T w (\\partial_x^{-1} h)^2\\,dx$, $h \\in H_S^\\bot$.\nRecalling \\eqref{def pi 0}, we calculate\n\\begin{align} \\label{L5 diff}\n& {\\cal L}_5 {\\cal S} - {\\cal S} \\Pi_S^\\bot\n( {\\cal D}_\\omega + m_3 \\partial_{xxx} + m_1 \\partial_x ) \\Pi_S^\\bot \n= \\Pi_S^\\bot ( 3 m_3 w_x + \\tilde{d}_1 - m_1 ) \\partial_x \\Pi_S^\\bot + \\tilde R_6 \\,,\n\\\\\n& \\tilde R_6\n:= \\Pi_S^\\bot \\{ ( 3 m_3 w_{xx} + \\tilde d_1 \\Pi_S^\\bot w - m_1 w ) \\pi_0 \n+ ( ({\\cal D}_\\omega w) + m_3 w_{xxx} + \\tilde d_1 \\Pi_S^\\bot w_x ) \\partial_x^{-1}\n+ ({\\cal D}_\\omega \\widehat{\\cal S}) \n\\notag \\\\ & \\qquad \\ \\ \n+ m_3 [\\partial_{xxx}, \\widehat{\\cal S}] \n+ \\tilde{d}_1 \\partial_x \\widehat{\\cal S} \n- m_1 \\widehat{\\cal S} \\partial_x \n+ R_5 {\\cal S} \\} \\Pi_S^\\bot\n\\notag \n\\end{align}\nwhere $\\tilde R_6$ collects all the terms of order at most $\\partial_x^0$. \nBy Remark \\ref{d1 media}, we solve $ 3 m_3 w_x + \\tilde{d}_1 - m_1 = 0 $ \nby choosing $w := - (3 m_3)^{-1} \\partial_x^{-1} ( \\tilde{d}_1 - m_1 )$.\nFor $ \\varepsilon $ small, the operator $ {\\cal S} $ is invertible and, by \\eqref{L5 diff}, \n\\begin{equation}\\label{def L6}\n\\mathcal{L}_6 := \\mathcal{S}^{-1} \\mathcal{L}_5 \\mathcal{S} \n= \\Pi_S^\\bot ( {\\cal D}_\\omega + m_3 \\partial_{xxx} + m_1 \\partial_x ) \\Pi_S^\\bot + R_6 \\,, \n\\qquad R_6 := {\\cal S}^{-1} \\tilde R_6 \\, .\n\\end{equation}\nSince $ {\\cal S} $ is symplectic, ${\\cal L}_6$ is Hamiltonian (recall Definition \\ref{operatore Hamiltoniano}).\n\n\n\\begin{lemma}\\label{lemma L6}\n There is $ \\sigma = \\sigma(\\nu,\\t) > 0 $ (possibly larger than in Lemma \\ref{lemma:R5}) such that \n$$\n|{\\cal S}^{\\pm 1} - I|_s^{\\mathrm{Lip}(\\g)} \n\\leq_s \\varepsilon^{5} \\gamma^{-1} + \\varepsilon \\| {\\mathfrak I}_\\delta\\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)}\\,, \\quad \n|\\partial_i {\\cal S}^{\\pm 1} [\\widehat \\imath ]|_s \n\\leq_s \\varepsilon ( \\| \\widehat \\imath\\|_{s + \\sigma} + \\|{\\mathfrak I}_\\delta \\|_{s + \\sigma} \\| \\widehat \\imath \\|_{s_0 + \\sigma} ).\n$$\nThe remainder $R_6$ satisfies the same estimates \\eqref{stima Lip R4} as $R_4$. \n\\end{lemma}\n\n\\begin{proof}\nBy \\eqref{tilde d1 d0 KdV},\\eqref{stima m1},\\eqref{stima m3}, $\\| w \\|_s^{\\mathrm{Lip}(\\g)} \\leq_s \\varepsilon^{5} \\gamma^{-1} + \\varepsilon \\| {\\mathfrak I}_\\delta\\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} $, and the lemma follows by \\eqref{def descent}. Since $ \\widehat { \\cal S } = O(\\partial_x^{-2})$ the commutator \n$ [ \\partial_{xxx}, \\widehat { \\cal S }] = O(\\partial_x^0 )$ and\n$ | [ \\partial_{xxx}, \\widehat { \\cal S }] |_s^{\\mathrm{Lip}(\\g)} \\leq_s \\| w \\|_{s_0+3}^{\\mathrm{Lip}(\\g)} \\| w \\|_{s+3}^{\\mathrm{Lip}(\\g)} $. \n\\end{proof}\n\n\\subsection{KAM reducibility and inversion of $ {\\cal L}_{\\omega} $} \\label{subsec:mL0 mL5}\n\nThe coefficients $ m_3, m_1 $ of the operator $ {\\cal L}_6 $ in \\eqref{def L6} are constants, and the remainder $ R_6 $ is a bounded operator of order $ \\partial_x^0 $ with small matrix decay \nnorm, see \\eqref{decay R6}.\nThen we can diagonalize $ {\\cal L}_6 $ \nby applying the iterative KAM reducibility Theorem 4.2 in \\cite{BBM} along the \nsequence of scales \n\\begin{equation}\\label{defN}\nN_n := N_{0}^{\\chi^n}, \\quad n = 0,1,2,\\ldots, \n\\quad \\chi := 3\/2, \\quad N_0 > 0 \\, .\n\\end{equation}\nIn section \\ref{sec:NM}, the initial $ N_0 $ will (slightly) increase to infinity \nas $ \\varepsilon \\to 0 $, see \\eqref{nash moser smallness condition}. \nThe required smallness condition (see (4.14) in \\cite{BBM}) is (written in the present notations)\n\\begin{equation}\\label{R6resto}\nN_0^{C_0} | R_6 |_{s_0 + \\beta}^{{\\mathrm{Lip}(\\g)}} \\gamma^{-1} \\leq 1 \n\\end{equation}\nwhere $ \\b := 7 \\tau + 6 $ (see (4.1) in \\cite{BBM}), \n$ \\tau $ is the diophantine exponent in \\eqref{omdio} and \\eqref{Omegainfty}, \nand the constant $ C_0 := C_0 (\\t, \\nu ) > 0 $ is fixed in\nTheorem 4.2 in \\cite{BBM}.\nBy Lemma \\ref{lemma L6}, the remainder $ R_6 $ satisfies the bound \\eqref{stima Lip R4}, \nand using \\eqref{ansatz delta} we get (recall \\eqref{link gamma b})\n\\begin{equation} \\label{decay R6}\n| R_6|_{s_0 + \\beta}^{{\\mathrm{Lip}(\\g)}} \\leq C \\varepsilon^{7 - 2 b} \\gamma^{-1} = C \\varepsilon^{3 - 2 a}, \\qquad\n| R_6 |_{s_0 + \\beta}^{{\\mathrm{Lip}(\\g)}} \\gamma^{-1} \\leq C \\varepsilon^{1 - 3 a} \\, .\n\\end{equation}\nWe use that $ \\mu $ in \\eqref{ansatz delta} is assumed to satisfy $ \\mu \\geq \\sigma + \\b $ where $ \\sigma := \\sigma (\\tau, \\nu ) $ \nis given in Lemma \\ref{lemma L6}. \n\n\\begin{theorem} \\label{teoremadiriducibilita} \n{\\bf (Reducibility)}\nAssume that $\\omega \\mapsto i_\\d (\\omega)\n$ is a Lipschitz function defined on some subset $\\Omega_o \\subset \\Omega_\\varepsilon $ (recall \\eqref{Omega epsilon}), satisfying \n\\eqref{ansatz delta} with $ \\mu \\geq \\sigma + \\b $ where $ \\sigma := \\sigma (\\tau, \\nu) $ is given in Lemma \\ref{lemma L6} and $ \\b := 7 \\tau + 6 $. \nThen there exists $ \\delta_{0} \\in (0,1) $ such that, if\n\\begin{equation}\\label{condizione-kam}\nN_0^{C_0} \\varepsilon^{7 - 2 b} \\gamma^{-2} = N_0^{C_0} \\varepsilon^{1 - 3 a}\n\\leq \\delta_{0} \\, , \\quad \\gamma := \\varepsilon^{2 + a} \\, , \\quad a \\in (0,1\/6) \\, , \n\\end{equation}\nthen:\n\n$(i)$ {\\bf (Eigenvalues)}.\nFor all $ \\omega \\in \\Omega_\\varepsilon $ there exists a sequence \n\\begin{equation} \\label{espressione autovalori}\n\\mu_j^\\infty(\\omega) := \\mu_j^\\infty(\\omega, i_\\d (\\omega)) \n:= {\\mathrm i} \\big( - {\\tilde m}_3 (\\omega) j^3 + {\\tilde m}_1(\\omega) j \\big)\n+ r_j^\\infty(\\omega), \\quad j \\in S^c \\, , \n\\end{equation}\nwhere $ {\\tilde m}_3, {\\tilde m}_1$ coincide with the coefficients $m_3, m_1$ of $ {\\cal L}_6 $ in \\eqref{def L6} for all $ \\omega \\in \\Omega_o $,\nand \n\\begin{align}\n\\label{autofinali}\n| {\\tilde m}_3 - 1 |^{{\\rm Lip}(\\gamma)} + | {\\tilde m}_1 |^{{\\rm Lip}(\\gamma)} \\leq C \\varepsilon^4 \\, , \\quad \n| r^{\\infty}_j |^{{\\rm Lip}(\\gamma)} & \\leq C \\varepsilon^{3 - 2 a} \\, , \n\\quad \\ \\forall j \\in S^c \\, , \n\\end{align}\nfor some $ C > 0 $.\nAll the eigenvalues $\\mu_j^{\\infty}$ are purely imaginary. We define, for convenience, $ \\mu_0^\\infty (\\omega) := 0 $. \n\n\\smallskip\n\n$(ii)$ {\\bf (Conjugacy)}.\nFor all $\\omega$ in the set\n\\begin{equation} \\label{Omegainfty}\n\\Omega_\\infty^{2\\gamma} := \\Omega_\\infty^{2\\gamma} (i_\\d) \n:= \\Big\\{ \\omega \\in \\Omega_o : \\,\n| {\\mathrm i} \\omega \\cdot l + \\mu^{\\infty}_j (\\omega) - \\mu^{\\infty}_{k} (\\omega) |\n\\geq \\frac{2 \\gamma | j^{3} - k^{3} |}{ \\langle l \\rangle^{\\tau}}, \n\\, \\forall l \\in \\mathbb Z^{\\nu}, \\, j ,k \\in S^c \\cup \\{0\\}\\Big\\} \n\\end{equation}\nthere is a real, bounded, invertible linear operator $\\Phi_\\infty(\\omega) : H^s_{S^\\bot} (\\mathbb T^{\\nu+1}) \\to \nH^s_{S^\\bot} (\\mathbb T^{\\nu+1}) $, with bounded inverse \n$\\Phi_\\infty^{-1}(\\omega)$, that conjugates $\\mathcal{L}_6$ in \\eqref{def L6} to constant coefficients, namely\n\\begin{equation}\\label{Lfinale}\n{\\cal L}_{\\infty}(\\omega)\n:= \\Phi_{\\infty}^{-1}(\\omega) \\circ \\mathcal{L}_6(\\omega) \\circ \\Phi_{\\infty}(\\omega)\n= \\omega \\cdot \\partial_{\\varphi} + {\\cal D}_{\\infty}(\\omega), \\quad\n{\\cal D}_{\\infty}(\\omega)\n:= {\\rm diag}_{j \\in S^c} \\{ \\mu^{\\infty}_{j}(\\omega) \\} \\, .\n\\end{equation}\nThe transformations $\\Phi_\\infty, \\Phi_\\infty^{-1}$ are close to the identity in matrix decay norm,\nwith\n\\begin{equation} \\label{stima Phi infty}\n| \\Phi_{\\infty} - I |_{s,\\Omega_\\infty^{2\\gamma}}^{{\\rm Lip}(\\gamma)}\n+ | \\Phi_{\\infty}^{- 1} - I |_{s,\\Omega_\\infty^{2\\gamma}}^{\\mathrm{Lip}(\\g)}\n\\leq_s \\varepsilon^{5} \\gamma^{-2} + \\varepsilon \\gamma^{-1} \\| {\\mathfrak I}_\\delta \\|_{s + \\sigma}^{\\mathrm{Lip}(\\g)} .\n\\end{equation}\nMoreover $\\Phi_{\\infty}, \\Phi_{\\infty}^{-1}$ are symplectic, and \n$\\mathcal{L}_\\infty $ is a Hamiltonian operator.\n\\end{theorem}\n\n\n\\begin{proof}\nThe proof is the same as the one of Theorem 4.1 in \\cite{BBM}, \nwhich is based on Theorem 4.2, Corollaries 4.1, 4.2 and Lemmata 4.1, 4.2 of \\cite{BBM}.\nA difference is that here $\\omega \\in \\mathbb R^\\nu$, while in \\cite{BBM} the parameter $\\lambda \\in \\mathbb R$ is one-dimensional. \nThe proof is the same because Kirszbraun's Theorem on Lipschitz extension of functions also holds in $\\mathbb R^\\nu$ \n(see, e.g., Lemma A.2 in \\cite{Po2}). \nThe bound \\eqref{stima Phi infty} follows by Corollary 4.1 of \\cite{BBM} and the estimate of \n$ R_6 $ in Lemma \\ref{lemma L6}.\nWe also use the estimates \\eqref{stima m3}, \\eqref{stima m1} for $ \\partial_i m_3 $, \n$ \\partial_i m_1 $ which correspond to (3.64) in \\cite{BBM}. \nAnother difference is that here the sites $ j \\in S^c \\subset \\mathbb Z \\setminus \\{0\\} $ unlike in \\cite{BBM} where $ j \\in \\mathbb Z $. \nWe have defined $ \\mu_0^\\infty := 0 $ \nso that also the first Melnikov conditions \\eqref{prime di melnikov} are included in the definition of $ \\Omega^{2 \\gamma}_{\\infty} $.\n\\end{proof}\n\n\\begin{remark}\nTheorem 4.2 in \\cite{BBM}\nalso provides the Lipschitz dependence of the (approximate) \neigenvalues $ \\mu_j^n $ with respect to the unknown $ i_0 (\\varphi) $, \nwhich is used for the measure estimate Lemma \\ref{matteo 10}.\n \\end{remark}\n\nAll the parameters $ \\omega \\in \\Omega_\\infty^{2 \\gamma} $ satisfy (specialize \\eqref{Omegainfty} for $ k = 0 $) \n\\begin{equation}\\label{prime di melnikov}\n|{\\mathrm i} \\omega \\cdot l + \\mu_j^\\infty(\\omega)| \\geq \n2 \\gamma | j |^3 \\langle l \\rangle^{-\\tau} \\, , \\quad \\forall l \\in \\mathbb Z^\\nu , \\ j \\in S^c, \n\\end{equation}\nand the diagonal operator $ {\\cal L}_\\infty $ is invertible. \n\nIn the following theorem we finally verify the inversion assumption \\eqref{tame inverse} for ${\\cal L}_\\omega $. \n\n\\begin{theorem}\\label{inversione linearized normale} {\\bf (Inversion of $ {\\cal L}_\\omega $)}\nAssume the hypotheses of Theorem \\ref{teoremadiriducibilita} and \\eqref{condizione-kam}. \nThen there exists $ \\sigma_1 := \\sigma_1 ( \\t, \\nu ) > 0 $ such that, \n$ \\forall \\omega \\in \\Omega^{2 \\gamma}_\\infty(i_\\d )$ (see \\eqref{Omegainfty}),\nfor any function $ g \\in H^{s+\\sigma_1}_{S^\\bot} (\\mathbb T^{\\nu+1}) $ \nthe equation ${\\cal L}_\\omega h = g$ \nhas a solution $h = {\\cal L}_\\omega^{-1} g \\in H^s_{S^\\bot} (\\mathbb T^{\\nu+1})$, satisfying\n\\begin{align}\\label{stima inverso linearizzato normale}\n\\| {\\cal L}_\\omega^{-1} g \\|_s^{{\\rm Lip}(\\gamma)} \n& \\leq_s \\gamma^{-1} \\big( \\| g \\|_{s +\\sigma_1}^{{\\rm Lip}(\\gamma)} \n+ \\varepsilon \\gamma^{-1} \\| {\\mathfrak I}_\\delta\\|_{s + \\sigma_1}^{\\mathrm{Lip}(\\g)}\n\\| g \\|_{s_0}^{{\\rm Lip}(\\gamma)} \\big) \n\\\\\n& \\leq_s \\gamma^{-1} \\big( \\| g \\|_{s +\\sigma_1}^{{\\rm Lip}(\\gamma)} \n+ \\varepsilon \\gamma^{-1} \\big\\{ \\| {\\mathfrak{I}}_0 \\|_{s + \\sigma_1 + \\sigma}^{\\mathrm{Lip}(\\g)} + \\gamma^{-1} \n \\| {\\mathfrak{I}}_0 \\|_{s_0 + \\sigma }^{\\mathrm{Lip}(\\g)}\n\\| Z \\|_{s + \\sigma_1 + \\sigma}^{\\mathrm{Lip}(\\g)} \\big\\}\n\\| g \\|_{s_0}^{{\\rm Lip}(\\gamma)} \\big)\\, . \\nonumber \n\\end{align}\n\\end{theorem}\n\n\\begin{proof}\nCollecting Theorem \\ref{teoremadiriducibilita} with the results of \nsections \\ref{step1}-\\ref{step5}, we have obtained the (semi)-conjugation of the operator \n$ \\mathcal{L}_\\omega $ (defined in \\eqref{Lom KdVnew}) \nto $\\mathcal{L}_\\infty $ (defined in \\eqref{Lfinale}), namely \n\\begin{equation} \\label{def cal M 12}\n\\mathcal{L}_\\omega = {\\cal M}_1 \\mathcal{L}_\\infty {\\cal M}_2^{-1}, \\qquad \n{\\cal M}_1 := \\Phi B \\rho {\\cal T} \\Phi_1 \\Phi_2 {\\cal S} \\Phi_{\\infty}, \\quad \n{\\cal M}_2 := \\Phi B {\\cal T} \\Phi_1 \\Phi_2 {\\cal S} \\Phi_{\\infty} \\,,\n\\end{equation}\nwhere $ \\rho $ means the multiplication operator by the function $ \\rho $ defined in \\eqref{anche def rho}. \nBy \\eqref{prime di melnikov} and Lemma 4.2 of \\cite{BBM} we deduce that \n$ \\| {\\cal L}_\\infty^{-1} g \\|_s^{\\mathrm{Lip}(\\g)} \\leq_s \\gamma^{-1} \\| g \\|_{s+ 2 \\tau + 1}^{\\mathrm{Lip}(\\g)} $.\nIn order to estimate $\\mathcal{M}_2, \\mathcal{M}_1^{-1}$, we recall that the composition of tame maps is tame, see Lemma 6.5 in \\cite{BBM}. \nNow, \n$ \\Phi , \\Phi^{-1}$ are estimated in Lemma \\ref{lemma:stime coeff mL1}, \n$ B, B^{-1} $ and $ \\rho $ in Lemma \\ref{lemma:stime coeff mL2}, \n$ {\\cal T}, {\\cal T}^{-1} $ in Lemma \\ref{lemma:stime coeff mL3}. \nThe decay norms $ |\\Phi_1|_s^{\\mathrm{Lip}(\\g)}$, $ |\\Phi_1^{-1}|_s^{\\mathrm{Lip}(\\g)} $, $ |\\Phi_2|_s^{\\mathrm{Lip}(\\g)} $, $ |\\Phi_2^{-1}|_s^{\\mathrm{Lip}(\\g)} \\leq C(s) $ by Lemmata \\ref{lemma:Dx A bounded}, \\ref{A2 decay}. \nThe decay norm of $ {\\cal S}, {\\cal S}^{-1} $ is estimated in Lemma \\ref{lemma L6}, \nand $ \\Phi_\\infty, \\Phi_\\infty^{-1} $ in \\eqref{stima Phi infty}. \nThe decay norm controls the Sobolev norm by \\eqref{interpolazione norme miste}. \nThus, by \\eqref{def cal M 12}, \n\\[\n\\| \\mathcal{M}_2 h \\|_s^{\\mathrm{Lip}(\\g)} + \\| \\mathcal{M}_1^{-1} h \\|_s^{\\mathrm{Lip}(\\g)} \n\\leq_s \\| h \\|_{s+3}^{\\mathrm{Lip}(\\g)} + \\varepsilon \\gamma^{-1} \\| {\\mathfrak{I}}_\\d \\|_{s+\\sigma+3}^{\\mathrm{Lip}(\\g)} \\| h \\|_{s_0}^{\\mathrm{Lip}(\\g)} \\,,\n\\]\nand \\eqref{stima inverso linearizzato normale} follows. \nThe last inequality in \\eqref{stima inverso linearizzato normale} follows by\n\\eqref{stima y - y delta} and \\eqref{ansatz 0}.\n\\end{proof}\n\n\n\\section{The Nash-Moser nonlinear iteration}\\label{sec:NM}\n \nIn this section we prove Theorem \\ref{main theorem}. It will be a consequence of the Nash-Moser Theorem \\ref{iterazione-non-lineare} below.\n\nConsider the finite-dimensional subspaces\n\\[\nE_n := \\big\\{ {\\mathfrak{I}} (\\varphi) = ( \\Theta, y, z )(\\varphi) : \\, \\Theta = \\Pi_n \\Theta, \\ y = \\Pi_n y, \\ z = \\Pi_n z \\big\\}\n\\]\nwhere $ N_n := N_0^{\\chi^n} $ are introduced in \\eqref{defN}, and $ \\Pi_n $ are the projectors \n(which, with a small abuse of notation, we denote with the same symbol)\n\\begin{align}\n\\Pi_n \\Theta (\\varphi) := \\sum_{|l| < N_n} \\Theta_l e^{{\\mathrm i} l \\cdot \\varphi}, \\quad \n\\Pi_n y (\\varphi) := \\sum_{|l| < N_n} y_l e^{{\\mathrm i} l \\cdot \\varphi}, \\quad \n& \\text{where} \\ \\Theta (\\varphi) = \\sum_{l \\in \\mathbb Z^\\nu} \\Theta_l e^{{\\mathrm i} l \\cdot \\varphi}, \\quad \ny(\\varphi) = \\sum_{l \\in \\mathbb Z^\\nu} y_l e^{{\\mathrm i} l \\cdot \\varphi}, \\nonumber\n\\\\\n\\Pi_n z(\\varphi,x) := \\sum_{|(l,j)| < N_n} z_{lj} e^{{\\mathrm i} (l \\cdot \\varphi + jx)}, \\quad \n& \\text{where} \\ z(\\varphi,x) = \\sum_{l \\in \\mathbb Z^\\nu, j \\in S^c} z_{lj} e^{{\\mathrm i} (l \\cdot \\varphi + jx)}. \\label{Pin def}\n\\end{align}\nWe define $ \\Pi_n^\\bot := I - \\Pi_n $. \nThe classical smoothing properties hold: for all $\\alpha , s \\geq 0$, \n\\begin{equation}\\label{smoothing-u1}\n\\|\\Pi_{n} {\\mathfrak{I}} \\|_{s + \\alpha}^{\\mathrm{Lip}(\\g)} \n\\leq N_{n}^{\\alpha} \\| {\\mathfrak{I}} \\|_{s}^{\\mathrm{Lip}(\\g)} \\, , \n\\ \\forall {\\mathfrak{I}} (\\omega) \\in H^{s} \\,, \n\\quad\n\\|\\Pi_{n}^\\bot {\\mathfrak{I}} \\|_{s}^{\\mathrm{Lip}(\\g)} \n\\leq N_{n}^{-\\alpha} \\| {\\mathfrak{I}} \\|_{s + \\alpha}^{\\mathrm{Lip}(\\g)} \\, , \n\\ \\forall {\\mathfrak{I}} (\\omega) \\in H^{s + \\alpha} \\, .\n\\end{equation}\nWe define the constants\n\\begin{alignat}{3} \\label{costanti nash moser}\n& \\mu_1 := 3 \\mu + 9\\,,\\quad &\n& \\alpha := 3 \\mu_1 + 1\\,,\\quad &\n& \\alpha_1 := (\\alpha - 3 \\mu)\/2 \\,, \n \\\\\n& \\kappa := 3 \\big(\\mu_1 + \\rho^{-1} \\big)+ 1\\,,\\qquad &\n& \\beta_1 := 6 \\mu_1+ 3 \\rho^{-1} + 3 \\, , \\qquad &\n& 0 < \\rho < \\frac{1 - 3 a}{C_1(1 + a)}\\,, \\label{def rho}\n\\end{alignat}\nwhere $ \\mu := \\mu (\\tau, \\nu) $ is the ``loss of regularity\" defined in Theorem \\ref{thm:stima inverso approssimato} \n(see \\eqref{stima inverso approssimato 1}) and $ C_1 $ is fixed below. \n\n\\begin{theorem}\\label{iterazione-non-lineare} \n{\\bf (Nash-Moser)} Assume that $ f \\in C^q $ with $ q > S := s_0 + \\b_1 + \\mu + 3 $. \nLet $ \\tau \\geq \\nu + 2 $. Then there exist $ C_1 > \\max \\{ \\mu_1 + \\a, C_0 \\} $\n(where $ C_0 := C_0 (\\tau, \\nu) $ is the one in Theorem \\ref{teoremadiriducibilita}), \n$ \\delta_0 := \\d_0 (\\tau, \\nu) > 0 $ such that, if\n\\begin{equation}\\label{nash moser smallness condition} \nN_0^{C_1} \\varepsilon^{b_* + 1} \\gamma^{-2}< \\d_0\\,,\\quad \\gamma:= \\varepsilon^{2 + a} = \\varepsilon^{2b} \\,,\\quad \nN_0 := (\\varepsilon \\gamma^{-1})^\\rho\\,,\\quad b_* := 6 - 2 b \\, , \n\\end{equation}\nthen, for all $ n \\geq 0 $: \n\n\\begin{itemize}\n\\item[$({\\cal P}1)_{n}$] \nthere exists a function \n$({\\mathfrak{I}}_n, \\zeta_n) : {\\cal G}_n \\subseteq \\Omega_\\varepsilon \\to E_{n-1} \\times \\mathbb R^\\nu$, \n$\\omega \\mapsto ({\\mathfrak{I}}_n(\\omega), \\zeta_n(\\omega))$, \n$ ({\\mathfrak{I}}_0, \\zeta_0) := 0 $, $ E_{-1} := \\{ 0 \\} $, \nsatisfying $ | \\zeta_n |^{\\mathrm{Lip}(\\g)} \\leq C \\|{\\cal F}(U_n) \\|_{s_0}^{\\mathrm{Lip}(\\g)} $, \n\\begin{equation}\\label{ansatz induttivi nell'iterazione}\n\\| {\\mathfrak{I}}_n \\|_{s_0 + \\mu }^{{\\rm Lip}(\\gamma)} \n\\leq C_* \\varepsilon^{b_*} \\gamma^{-1}\\,, \\quad \n\\| {\\cal F}(U_n)\\|_{s_0 + \\mu + 3}^{{\\rm Lip}(\\gamma)} \\leq C_*\\varepsilon^{b_*} \\,, \n\\end{equation}\nwhere $U_n := (i_n, \\zeta_n)$ with $i_n(\\varphi) = (\\varphi,0,0) + {\\mathfrak{I}}_n(\\varphi)$.\nThe sets ${\\cal G}_{n} $ are defined inductively by: \n$$\n{\\cal G}_{0} := \\big\\{\\omega \\in \\Omega_\\varepsilon \\, :\\, |\\omega \\cdot l| \\geq 2 \\gamma \\langle l \\rangle^{-\\tau}, \\, \\forall l \\in \\mathbb Z^\\nu \\setminus \\{0\\} \\big\\}\\,, \n$$\n\\begin{equation}\\label{def:Gn+1}\n{\\cal G}_{n+1} := \n\\Big\\{ \\omega \\in {\\cal G}_{n} \\, : \\, |{\\mathrm i} \\omega \\cdot l + \\mu_j^\\infty ( i_n) -\n\\mu_k^\\infty ( i_n )| \\geq \\frac{2\\gamma_{n} |j^{3}-k^{3}|}{\\left\\langle l\\right\\rangle^{\\tau}}, \\, \\forall j , k \\in S^c \\cup \\{0\\}, \n\\, l \\in \\mathbb Z^{\\nu} \\Big\\}\\,,\n\\end{equation}\nwhere $ \\gamma_{n}:=\\gamma (1 + 2^{-n}) $ and $\\mu_j^\\infty(\\omega) := \\mu_j^\\infty(\\omega, i_n(\\omega)) $ \nare defined in \\eqref{espressione autovalori} (and $ \\mu_0^\\infty(\\omega) = 0 $).\n\nThe differences $\n\\widehat {\\mathfrak I}_n\n:= {\\mathfrak I}_n - {\\mathfrak I}_{n - 1} $\n(where we set $ \\widehat {\\mathfrak{I}}_0 := 0 $) is defined on $\\mathcal{G}_n$, and satisfy\n\\begin{equation} \\label{Hn}\n\\| \\widehat {\\mathfrak I}_1 \\|_{ s_0 + \\mu}^{{\\mathrm{Lip}(\\g)}} \\leq C_* \\varepsilon^{b_*} \\gamma^{-1} \\, , \\quad \n\\| \\widehat {\\mathfrak I}_n \\|_{ s_0 + \\mu}^{{\\mathrm{Lip}(\\g)}} \\leq C_* \\varepsilon^{b_*} \\gamma^{-1} N_{n - 1}^{-\\alpha_1} \\, , \\quad \\forall n > 1 \\, .\n\\end{equation}\n\\item[$({\\cal P}2)_{n}$] $ \\| {\\cal F}(U_n) \\|_{ s_{0}}^{{\\rm Lip}(\\gamma)} \\leq C_* \\varepsilon^{b_*} N_{n - 1}^{- \\alpha}$ \nwhere we set $N_{-1} := 1$.\n\\item[$({\\cal P}3)_{n}$] \\emph{(High norms).} \n\\ $ \\| {\\mathfrak{I}}_n \\|_{ s_{0}+ \\beta_1}^{{\\rm Lip}(\\gamma)} \\leq C_* \\varepsilon^{b_*} \\gamma^{-1} N_{n - 1}^{\\kappa} $ and \n$ \\|{\\cal F}(U_n ) \\|_{ s_{0}+\\beta_1}^{{\\rm Lip}(\\gamma)} \\leq C_* \\varepsilon^{b_*} N_{n - 1}^{\\kappa} $. \n\n\n\\item[$({\\cal P}4)_{n}$] \\emph{(Measure).} \nThe measure of the ``Cantor-like\" sets $ {\\cal G}_n $ satisfies\n\\begin{equation}\\label{Gmeasure}\n | \\Omega_\\varepsilon \\setminus {\\cal G}_0 | \\leq C_* \\varepsilon^{2(\\nu - 1)} \\gamma \\, , \\quad \n\\big| {\\cal G}_n \\setminus {\\cal G}_{n+1} \\big| \\leq C_* \\varepsilon^{2(\\nu - 1)} \\gamma N_{n - 1}^{-1} \\, . \n\\end{equation}\n\\end{itemize}\nAll the Lip norms are defined on $ {\\cal G}_{n} $, namely $\\| \\ \\|_s^{{\\rm Lip}(\\gamma)} = \\| \\ \\|_{s,\\mathcal{G}_n}^{{\\rm Lip}(\\gamma)} $.\\end{theorem}\n\n\n\\begin{proof}\nTo simplify notations, in this proof we denote $\\| \\, \\|^{{\\rm Lip}(\\gamma)}$ by $\\| \\, \\|$. We first prove $({\\cal P}1, 2, 3)_n$.\n\n\\smallskip\n\n{\\sc Step 1:} \\emph{Proof of} $({\\cal P}1, 2, 3)_0$. \nRecalling \\eqref{operatorF} we have $ \\| {\\cal F}( U_0 ) \\|_s = $ \n$ \\| {\\cal F}(\\varphi, 0 , 0, 0 ) \\|_s = \\| X_P(\\varphi, 0 , 0 ) \\|_s \\leq_s \\varepsilon^{6 - 2b} $ by\n\\eqref{stima XP}. Hence (recall that $ b_* = 6 - 2 b $) the smallness conditions in\n$({\\cal P}1)_0$-$({\\cal P}3)_0$ hold taking $ C_* := C_* (s_0 + \\b_1) $ large enough.\n\n\\smallskip\n\n{\\sc Step 2:} \\emph{Assume that $({\\cal P}1,2,3)_n$ hold for some $n \\geq 0$, and prove $({\\cal P}1,2,3)_{n+1}$.}\nBy \\eqref{nash moser smallness condition} and \\eqref{def rho},\n$$\nN_0^{C_1} \\varepsilon^{b_* + 1} \\gamma^{-2} = N_0^{C_1} \\varepsilon^{1-3a} = \\varepsilon^{1 - 3 a - \\rho C_1(1 + a)} < \\d_0 \n$$\nfor $\\varepsilon$ small enough, and the smallness condition \\eqref{condizione-kam} holds. \nMoreover \\eqref{ansatz induttivi nell'iterazione} imply \\eqref{ansatz 0}\n(and so \\eqref{ansatz delta}) and Theorem \\ref{inversione linearized normale} applies. \nHence the operator $ {\\cal L}_\\omega := {\\cal L}_\\omega(\\omega, i_n(\\omega))$ \ndefined in \\eqref{cal L omega} is invertible for all $\\omega \\in {\\cal G}_{n + 1}$ and \nthe last estimate in \\eqref{stima inverso linearizzato normale} holds. \nThis means that the assumption \\eqref{tame inverse} of Theorem \\ref{thm:stima inverso approssimato} \nis verified with $\\Omega_\\infty = {\\cal G}_{n + 1}$. \nBy Theorem \\ref{thm:stima inverso approssimato} there exists an approximate inverse \n${\\bf T}_n(\\omega) := {\\bf T}_0 (\\omega, i_n(\\omega))$ of the linearized operator \n$L_n(\\omega) := d_{i, \\zeta} {\\cal F}(\\omega, i_n(\\omega)) $, satisfying \\eqref{stima inverso approssimato 1}. \nThus, using also \\eqref{nash moser smallness condition}, \\eqref{smoothing-u1}, \\eqref{ansatz induttivi nell'iterazione},\n\\begin{align}\n\\| {\\bf T}_n g \\|_s \n& \\leq_s \\gamma^{-1} \\big( \\| g \\|_{s + \\mu} \n+ \\varepsilon \\gamma^{-1}\\{\\| {\\mathfrak{I}}_n \\|_{s + \\mu} + \n\\gamma^{-1} \\|{\\mathfrak I}_n \\|_{s_0 + \\mu} \\|{\\cal F}(U_n) \\|_{s + \\mu} \\} \\| g \\|_{s_0+ \\mu}\\big) \\label{stima Tn} \\\\\n\\| {\\bf T}_n g \\|_{s_0} \n& \n\\leq_{s_0} \\gamma^{-1} \\| g \\|_{s_0 + \\mu} \\label{stima Tn norma bassa}\n\\end{align}\nand, by \\eqref{stima inverso approssimato 2}, \nusing also \\eqref{ansatz induttivi nell'iterazione}, \n\\eqref{nash moser smallness condition}, \\eqref{smoothing-u1}, \n\\begin{align}\n\\| \\big(L_n \\circ {\\bf T}_n -I \\big) g \\|_s \n& \\leq_s \\gamma^{-1} \\big( \\| {\\cal F}(U_n) \\|_{s_0 + \\mu} \\| g\\|_{s + \\mu} + \n\\| {\\cal F}(U_n) \\|_{s + \\mu} \\| g \\|_{s_0 + \\mu} \\nonumber\n\\\\\n& \\qquad + \\varepsilon \\gamma^{-1} \\| {\\mathfrak{I}}_n \\|_{s + \\mu} \\| {\\cal F}(U_n) \\|_{s_0 + \\mu} \\| g \\|_{s_0 + \\mu} \\big)\\, , \n \\label{stima Tn inverso approssimato} \\\\\n\\label{stima Tn inverso approssimato norma bassa}\n\\| \\big(L_n \\circ {\\bf T}_n -I \\big) g \\|_{s_0} \n& \\leq_{s_0} \\gamma^{-1} \\| {\\cal F}(U_n)\\|_{s_0 + \\mu} \\| g \\|_{s_0 + \\mu} \\nonumber\n\\\\\n& \\leq_{s_0} \\gamma^{-1} \\big( \\|\\Pi_n {\\cal F}(U_n)\\|_{s_0 + \\mu}+ \\|\\Pi_n^\\bot {\\cal F}(U_n)\\|_{s_0 + \\mu} \\big) \\| g \\|_{s_0 + \\mu} \\nonumber \n\\\\\n& \\leq_{s_0} N_{n }^{\\mu} \\gamma^{-1} \\big( \\| {\\cal F}(U_n)\\|_{s_0} + N_{n}^{-\\beta_1} \\| {\\cal F}(U_n)\\|_{s_0 + \\beta_1} \\big) \\| g \\|_{s_0 + \\mu}\\, .\n\\end{align}\nThen, for all $ \\omega \\in {\\cal G}_{n+1} $, $ n \\geq 0 $, we define \n\\begin{equation}\\label{soluzioni approssimate}\nU_{n + 1} := U_n + H_{n + 1}\\,, \\quad \nH_{n + 1} :=\n( \\widehat {\\mathfrak{I}}_{n+1}, \\widehat \\zeta_{n+1}) := - {\\widetilde \\Pi}_{n } {\\bf T}_n \\Pi_{n } {\\cal F}(U_n) \n\\in E_n \\times \\mathbb R^\\nu\\, , \n\\end{equation}\nwhere $ {\\widetilde \\Pi}_n ( {\\mathfrak{I}} , \\zeta ) := ( \\Pi_n {\\mathfrak{I}} , \\zeta ) $ \nwith $ \\Pi_n $ in \\eqref{Pin def}.\nSince $ L_n := d_{i,\\zeta} {\\cal F}(i_n) $, we write \n$ {\\cal F}(U_{n + 1}) = {\\cal F}(U_n) + L_n H_{n + 1} + Q_n $, where\n\\begin{equation}\\label{def:Qn}\nQ_n := Q(U_n, H_{n + 1}) \\, , \\quad \nQ (U_n, H) := {\\cal F}(U_n + H ) - {\\cal F}(U_n) - L_n H \\,, \\quad \nH \\in E_{n} \\times \\mathbb R^\\nu . \n\\end{equation}\nThen, by the definition of $ H_{n+1} $ in \\eqref{soluzioni approssimate}, and writing $ {\\widetilde \\Pi}_n^\\bot ({\\mathfrak{I}}, \\zeta ) :=\n(\\Pi_n^\\bot {\\mathfrak{I}}, 0) $, we have \n\\begin{align}\n{\\cal F}(U_{n + 1}) & = \n {\\cal F}(U_n) - L_n {\\widetilde \\Pi}_{n } {\\bf T}_n \\Pi_{n } {\\cal F}(U_n) + Q_n = \n {\\cal F}(U_n) - L_n {\\bf T}_n \\Pi_{n } {\\cal F}(U_n) + L_n {\\widetilde \\Pi}_n^\\bot {\\bf T}_n \\Pi_{n } {\\cal F}(U_n)\n + Q_n \\nonumber\\\\ \n& = {\\cal F}(U_n) - \\Pi_{n } L_n {\\bf T}_n \\Pi_{n }{\\cal F}(U_n) \n+ ( L_n {\\widetilde \\Pi}_n^\\bot - \\Pi_n^\\bot L_n ) {\\bf T}_n \\Pi_{n }{\\cal F}(U_n) + Q_n \\nonumber\\\\\n & = \\Pi_{n }^\\bot {\\cal F}(U_n) + R_n + Q_n + Q_n' \n\\label{relazione algebrica induttiva}\n\\end{align}\nwhere \n\\begin{equation}\\label{Rn Q tilde n}\nR_n := (L_n {\\widetilde \\Pi}_n^\\bot - \\Pi_n^\\bot L_n) {\\bf T}_n \\Pi_{n }{\\cal F}(U_n) \\,,\n\\qquad \nQ_n' := - \\Pi_{n } ( L_n {\\bf T}_n - I) \\Pi_{n } {\\cal F}(U_n)\\,.\n\\end{equation}\n\n\n\\begin{lemma}\\label{lemma convergence}\nDefine \n\\begin{equation}\\label{riscalamenti nash moser}\nw_n := \\varepsilon \\gamma^{-2} \\|{\\cal F}(U_n) \\|_{s_0}\\,,\\quad B_n := \\varepsilon \\gamma^{-1}\\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1} + \\varepsilon \\gamma^{-2} \\|{\\cal F}(U_n) \\|_{s_0 + \\beta_1} \\,.\n\\end{equation}\nThen there exists $ K := K( s_0, \\b_1 ) > 0 $ such that, for all $n \\geq 0$, \nsetting $ \\mu_1 := 3 \\mu + 9$ (see \\eqref{costanti nash moser}), \n\\begin{equation}\\label{relazioni induttive}\nw_{n + 1} \\leq K N_{n }^{\\mu_1 + \\frac{1}{\\rho} - \\beta_1} B_n + K N_n^{\\mu_1} w_n^2\\,,\\qquad \nB_{n + 1} \\leq K N_{n }^{\\mu_1 + \\frac{1}{\\rho}} B_n\\, . \n\\end{equation}\n\\end{lemma}\n\\begin{proof}\nWe estimate separately the terms $ Q_n $ in \\eqref{def:Qn} and $ Q_n' , R_n $ in \\eqref{Rn Q tilde n}. \n\\\\[1mm]\n{\\it Estimate of $ Q_n $.}\nBy \\eqref{def:Qn}, \\eqref{operatorF}, \\eqref{parte quadratica da P}\nand \\eqref{ansatz induttivi nell'iterazione}, \\eqref{smoothing-u1}, \nwe have the quadratic estimates\n\\begin{align}\n\\label{stima parte quadratica norma alta}\n\\| Q(U_n, H)\\|_s & \\leq_s \\varepsilon \\big(\\| \\widehat {\\mathfrak{I}} \\|_{s + 3} \\| \\widehat {\\mathfrak{I}} \\|_{s_0 + 3} + \\| {\\mathfrak{I}}_n \\|_{s + 3} \n\\| \\widehat {\\mathfrak{I}} \\|_{s_0 + 3}^2 \\big) \\\\\n\\label{stima parte quadratica norma bassa}\n\\| Q(U_n, H) \\|_{s_0} & \\leq_{s_0}\\varepsilon N_n^6 \\| \\widehat {\\mathfrak{I}} \\|_{s_0}^2\\,, \\quad \\forall \\widehat {\\mathfrak{I}} \\in E_n \\, . \n\\end{align}\nNow by the definition of $H_{n + 1}$ in \\eqref{soluzioni approssimate} and \\eqref{smoothing-u1}, \n\\eqref{stima Tn}, \\eqref{stima Tn norma bassa}, \\eqref{ansatz induttivi nell'iterazione}, we get \n\\begin{align}\n\\| \\widehat {\\mathfrak{I}}_{n + 1} \\|_{s_0 + \\beta_1} \n& \\leq_{s_0 + \\beta_1} N_n^{\\mu} \\big( \\gamma^{-1} \\|{\\cal F}(U_n) \\|_{s_0 + \\beta_1} + \\varepsilon \\gamma^{-2} \\| {\\cal F}(U_n)\\|_{s_0 + \\mu} \\{\\| {\\mathfrak{I}}_n\\|_{s_0 + \\beta_1} + \\gamma^{-1} \\| {\\cal F}(U_n)\\|_{s_0 + \\beta_1} \\} \\big)\n\\nonumber \\\\\n& \\leq_{s_0 + \\beta} N_n^{\\mu} \\big( \\gamma^{-1}\\| {\\cal F}(U_n) \\|_{s_0 + \\beta_1} \n+ \\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1}\\big)\\, , \\label{H n+1 alta} \n\\\\\n\\label{H n+1 bassa}\n\\| \\widehat {\\mathfrak{I}}_{n + 1}\\|_{s_0} \n& \\leq_{s_0} \\gamma^{-1}N_{n}^\\mu \\| {\\cal F}(U_n)\\|_{s_0} \\, .\n\\end{align}\nThen the term $ Q_n $ in \\eqref{def:Qn} satisfies, \nby \\eqref{stima parte quadratica norma alta}, \\eqref{stima parte quadratica norma bassa}, \n\\eqref{H n+1 alta}, \\eqref{H n+1 bassa}, \n\\eqref{nash moser smallness condition}, \\eqref{ansatz induttivi nell'iterazione}, \n$ ({\\cal P}2)_n $, \\eqref{costanti nash moser}, \n\\begin{align} \n\\| Q_n \\|_{s_0 + \\beta_1} \n& \\leq_{s_0 + \\beta_1}\n N_n^{2 \\mu + 9} \\gamma \\big( \\gamma^{-1} \\| {\\cal F}(U_n)\\|_{s_0 + \\beta_1} + \n\\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1} \\big) \\, , \\label{Qn norma alta} \\\\\n\\| Q_n \\|_{s_0} \n& \\leq_{s_0} N_n^{2 \\mu + 6 } \\varepsilon \\gamma^{-2} \\| {\\cal F}(U_n) \\|_{s_0}^2\\, . \\label{Qn norma bassa}\n\\end{align}\n{\\it Estimate of $ Q_n' $.}\nThe bounds \\eqref{stima Tn inverso approssimato}, \\eqref{stima Tn inverso approssimato norma bassa},\n\\eqref{smoothing-u1}, \\eqref{costanti nash moser}, \\eqref{ansatz induttivi nell'iterazione} imply\n\\begin{align}\n\\label{Qn' norma alta}\n\\| Q_n' \\|_{s_0 + \\beta_1} & \\leq_{s_0 + \\beta_1} N_n^{2 \\mu} \\big(\n\\| {\\cal F}(U_n) \\|_{s_0 + \\beta_1} + \\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1} \\| {\\cal F}(U_n)\\|_{s_0} \\big) \\, , \\\\\n\\label{Qn' norma bassa}\n\\| Q_n'\\|_{s_0} & \\leq_{s_0} \\gamma^{-1} N_{n }^{2\\mu }\\big(\\| {\\cal F}(U_n) \\|_{s_0} + N_{n}^{- \\beta_1} \\|{\\cal F}(U_n)\\|_{s_0 + \\beta_1} \\big) \\| {\\cal F}(U_n) \\|_{s_0}\\,.\n\\end{align}\n{\\it Estimate of $ R_n $.} For $ H := (\\widehat {\\mathfrak{I}}, \\widehat \\zeta ) $ we have \n$ (L_n {\\widetilde \\Pi}_n^\\bot - \\Pi_n^\\bot L_n) H = $ $ [{\\bar D}_n, \\Pi_n^\\bot ] \\widehat {\\mathfrak{I}} = $\n$ [\\Pi_n , {\\bar D}_n] \\widehat {\\mathfrak{I}} $ where $ {\\bar D}_n := d_i X_{H_\\varepsilon}(i_n) + (0,0, \\partial_{xxx} ) $. \nThus Lemma \\ref{lemma quantitativo forma normale},\n\\eqref{ansatz induttivi nell'iterazione}, \\eqref{smoothing-u1} and \\eqref{tame commutatori} imply \n\\begin{align}\\label{stima commutatore modi alti norma bassa}\n\\| (L_n {\\widetilde \\Pi}_n^\\bot - \\Pi_n^\\bot L_n) H \\|_{s_0} & \\leq_{s_0+ \\b_1} \n\\varepsilon N_{n }^{- \\beta_1 + \\mu + 3} \\big(\\| \\widehat {\\mathfrak{I}} \\|_{s_0 + \\beta_1 - \\mu} + \n\\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1 - \\mu} \\| \\widehat {\\mathfrak{I}} \\|_{s_0 + 3}\\big)\\,, \\\\\n\\label{stima commutatore modi alti norma alta}\n \\| (L_n {\\widetilde \\Pi}_n^\\bot - \\Pi_n^\\bot L_n) H \\|_{s_0 + \\beta_1} & \n\\leq_s \n\\varepsilon N_n^{\\mu + 3} \\big(\\| \\widehat {\\mathfrak{I}} \\|_{s_0 + \\beta_1 - \\mu } + \\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1 - \\mu } \n\\| \\widehat {\\mathfrak{I}} \\|_{s_0 + 3} \\big)\\,.\n\\end{align}\nHence, applying \n\\eqref{stima Tn}, \n\\eqref{stima commutatore modi alti norma bassa}, \\eqref{stima commutatore modi alti norma alta}, \n\\eqref{nash moser smallness condition}, \n\\eqref{ansatz induttivi nell'iterazione}, \n\\eqref{smoothing-u1}, \nthe term $R_n$ defined in \\eqref{Rn Q tilde n} satisfies\n\\begin{align} \\label{stima Rn norma bassa}\n\\| R_n\\|_{s_0} \n& \\leq_{s_0 + \\beta_1} N_n^{ \\mu + 6 - \\beta_1} ( \\varepsilon \\gamma^{-1} \\| {\\cal F}(U_n)\\|_{s_0 + \\beta_1} + \\varepsilon \\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1} )\\,, \n\\\\\n\\label{stima Rn norma alta}\n\\| R_n \\|_{s_0 + \\beta_1} \n& \\leq_{s_0 + \\beta_1} N_n^{ \\mu + 6} ( \\varepsilon \\gamma^{-1} \\| {\\cal F}(U_n)\\|_{s_0 + \\beta_1} + \\varepsilon \\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1} )\\,.\n\\end{align}\n{\\it Estimate of $ {\\cal F}(U_{n + 1}) $.} \nBy \\eqref{relazione algebrica induttiva} and \\eqref{Qn norma alta}, \n \\eqref{Qn norma bassa}, \\eqref{Qn' norma alta}, \\eqref{Qn' norma bassa}, \n\\eqref{stima Rn norma bassa}, \\eqref{stima Rn norma alta}, \\eqref{nash moser smallness condition}, \n\\eqref{ansatz induttivi nell'iterazione}, \nwe get\n\\begin{align}\\label{F(U n+1) norma bassa}\n& \\| {\\cal F}(U_{n + 1})\\|_{s_0} \\leq_{s_0 + \\beta_1} N_{n }^{\\mu_1 - \\beta_1} ( \\varepsilon \\gamma^{-1}\\| {\\cal F}(U_n)\\|_{s_0 + \\beta_1} + \\varepsilon \\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1} ) \n+ N_n^{\\mu_1} \\varepsilon \\gamma^{-2} \\| {\\cal F}(U_n)\\|_{s_0}^2\\,, \n\\\\\n\\label{F(U n+1) norma alta}\n& \\| {\\cal F}(U_{n + 1}) \\|_{s_0 + \\beta_1} \n\\leq_{s_0 + \\beta_1} N_n^{\\mu_1} ( \\varepsilon \\gamma^{-1}\\| {\\cal F}(U_n) \\|_{s_0 + \\beta_1} + \\varepsilon \\| {\\mathfrak{I}}_n \\|_{s_0 + \\beta_1} ) \\,,\n\\end{align}\nwhere $\\mu_1 := 3 \\mu + 9$.\n\n\\noindent \n{\\it Estimate of $ {\\mathfrak{I}}_{n+1} $.} \nUsing \\eqref{H n+1 alta} the term $ {\\mathfrak{I}}_{n+1} = {\\mathfrak{I}}_n + \\widehat {\\mathfrak I}_{n+1} $ is bounded by\n\\begin{equation}\\label{U n+1 alta}\n\\| {\\mathfrak{I}}_{n + 1}\\|_{s_0 + \\beta_1} \\leq_{s_0 + \\beta_1} \nN_n^\\mu ( \\| {\\mathfrak{I}}_n\\|_{s_0 + \\beta_1} +\\gamma^{-1} \\| {\\cal F}(U_n)\\|_{s_0 + \\beta_1} )\\, . \n\\end{equation}\nFinally, recalling \\eqref{riscalamenti nash moser}, the inequalities \\eqref{relazioni induttive} follow by \n\\eqref{F(U n+1) norma bassa}-\\eqref{U n+1 alta}, \\eqref{ansatz induttivi nell'iterazione} and \n$ \\varepsilon \\gamma^{-1} = N_0^{1\/\\rho} \\leq N_n^{1\/\\rho}$.\n\\end{proof} \n\n\\emph{Proof of $({\\cal P}3)_{n + 1}$}. \nBy \\eqref{relazioni induttive} and $ ({\\cal P}3)_n $, \n\\begin{equation}\nB_{n + 1} \n\\leq K N_{n}^{\\mu_1 + \\frac{1}{\\rho}} B_n \n\\leq 2 C_* K \\varepsilon^{b_* + 1} \\gamma^{-2} N_{n}^{\\mu_1+ \\frac{1}{\\rho}} N_{n-1}^\\kappa \\leq C_* \\varepsilon^{b_* + 1} \\gamma^{-2} N_n^{\\kappa} \\,, \n\\label{stima B n+1} \n\\end{equation}\nprovided $ 2 K N_n^{\\mu_1 + \\frac{1}{\\rho} - \\kappa} N_{n-1}^\\kappa \\leq 1 $, $ \\forall n \\geq 0 $. \nThis inequality holds by \\eqref{def rho}, taking $N_0$ large enough (i.e $\\varepsilon$ small enough). By \\eqref{riscalamenti nash moser}, the bound $B_{n + 1} \\leq C_* \\varepsilon^{b_* + 1} \\gamma^{-2} N_n^{\\kappa}$ implies $({\\cal P}3)_{n + 1}$.\n\n\\smallskip\n\n\\emph{Proof of $({\\cal P}2)_{n + 1}$}. \nUsing \\eqref{relazioni induttive}, \\eqref{riscalamenti nash moser} and $ ({\\cal P}2)_n, ({\\cal P}3)_n $, we get\n\\begin{align*}\nw_{n + 1} \n& \\leq K N_n^{\\mu_1+ \\frac{1}{\\rho} - \\beta_1} B_n + K N_{n }^{\\mu_1} w_n^2 \n\\leq K N_n^{\\mu_1 + \\frac{1}{\\rho} - \\beta_1} 2 C_* \\varepsilon^{b_* + 1} \\gamma^{-2} N_{n - 1}^{\\kappa} \n+ K N_{n}^{\\mu_1} ( C_* \\varepsilon^{b_* + 1} \\gamma^{-2} N_{n - 1}^{- \\alpha} )^2 \n\\end{align*}\nwhich is $ \\leq C_* \\varepsilon^{b_* + 1} \\gamma^{-2} N_n^{- \\alpha} $ provided that \n\\begin{equation} \\label{provided 2}\n4 K N_n^{\\mu_1 + \\frac{1}{\\rho} - \\b_1 + \\a} N_{n-1}^\\kappa \\leq 1, \n\\quad\n2 K C_* \\varepsilon^{b_*+1} \\gamma^{-2} N_n^{\\mu_1 + \\a} N_{n-1}^{-2\\a} \\leq 1 \\, , \n\\quad \\forall n \\geq 0.\n\\end{equation}\nThe inequalities in \\eqref{provided 2} hold by \\eqref{costanti nash moser}-\\eqref{def rho}, \n\\eqref{nash moser smallness condition}, $C_1 > \\mu_1 + \\alpha$, taking $\\d_0$ in \\eqref{nash moser smallness condition} small enough. \nBy \\eqref{riscalamenti nash moser}, the inequality $w_{n + 1} \\leq C_* \\varepsilon^{b_* + 1} \\gamma^{-2} N_n^{- \\alpha}$ implies $({\\cal P}2)_{n + 1}$.\n\n\\smallskip\n\n\\emph{Proof of $({\\cal P}1)_{n + 1}$}.\nThe bound \\eqref{Hn} for $ \\widehat {\\mathfrak{I}}_1 $ \nfollows by \\eqref{soluzioni approssimate}, \\eqref{stima Tn} (for $ s = s_0 + \\mu $) and \n$ \\| {\\cal F} ( U_0 ) \\|_{s_0 + 2 \\mu} = $ $ \\| {\\cal F}(\\varphi, 0, 0, 0) \\|_{s_0 + 2 \\mu} \\leq_{s_0+2\\mu } \\varepsilon^{b_*} $.\nThe bound \\eqref{Hn} for $ \\widehat {\\mathfrak{I}}_{n + 1} $ follows by \\eqref{smoothing-u1}, \n\\eqref{H n+1 bassa}, $({\\cal P}2)_n $, \\eqref{costanti nash moser}. \nIt remains to prove that \\eqref{ansatz induttivi nell'iterazione} holds at the step $n + 1$.\nWe have\n\\begin{equation}\\label{W n+1 norma bassa}\n\\| {\\mathfrak{I}}_{n + 1} \\|_{s_0 + \\mu} \n\\leq {\\mathop \\sum}_{k = 1}^{n + 1} \\| \\widehat {\\mathfrak I}_k \\|_{s_0 + \\mu} \n\\leq C_* \\varepsilon^{b_*}\\gamma^{-1} {\\mathop \\sum}_{k \\geq 1} N_{k - 1}^{- \\alpha_1} \n\\leq C_* \\varepsilon^{b_*} \\gamma^{-1}\n\\end{equation}\nfor $ N_0 $ large enough, i.e. $ \\varepsilon $ small. \nMoreover, using \\eqref{smoothing-u1}, (${\\cal P}2)_{n +1}$, (${\\cal P}3)_{n + 1}$, \\eqref{costanti nash moser}, we get \n\\begin{align*}\n\\| {\\cal F}(U_{n + 1}) \\|_{s_0 + \\mu+ 3} \n& \\leq N_n^{\\mu + 3} \\|{\\cal F}(U_{n + 1}) \\|_{s_0} + N_{n}^{\\mu + 3 - \\beta_1} \\| {\\cal F}(U_{n + 1}) \\|_{s_0 + \\beta_1} \n\\\\\n& \\leq C_* \\varepsilon^{b_*} N_n^{\\mu + 3- \\alpha} + C_* \\varepsilon^{b_*} N_n^{\\mu + 3 - \\beta_1 + \\kappa}\n\\leq C_* \\varepsilon^{b_*}\\,,\n\\end{align*}\nwhich is the second inequality in \\eqref{ansatz induttivi nell'iterazione} at the step $n + 1$.\nThe bound $ | \\zeta_{n+1} |^{\\mathrm{Lip}(\\g)} \\leq C \\|{\\cal F}(U_{n+1}) \\|_{s_0}^{\\mathrm{Lip}(\\g)} $ is a consequence of Lemma \n\\ref{zeta = 0} (it is not inductive).\n\n\\smallskip\n\n\n{\\sc Step 3:} \\emph{Prove $({\\cal P}4)_n$ for all $n \\geq 0$.} \nFor all $n \\geq 0$, \n\\begin{equation}\\label{Gn inclusioni}\n{\\cal G}_n \\setminus {\\cal G}_{n + 1} = \\!\\! \\!\\!\\! \\bigcup_{\\begin{subarray}{c}\nl \\in \\mathbb Z^\\nu, \\, j , k \\in S^c \\cup\\{ 0 \\}\n\\end{subarray}} \\!\\! \\! \\!\\! R_{ljk}(i_n)\n\\end{equation}\nwhere \n\\begin{equation}\\label{resonant sets}\nR_{ljk}(i_n) := \\big\\{ \\omega \\in {\\cal G}_n \\, : \\, |{\\mathrm i} \\omega \\cdot l + \\mu_j^\\infty (i_{n}) -\n\\mu_k^\\infty (i_{n})| < 2\\gamma_{n} |j^{3}-k^{3}|\\left\\langle l\\right\\rangle^{- \\tau}\\big\\}\\,.\n\\end{equation}\nNotice that $R_{ljk}(i_n) = \\emptyset$ if $j = k$, so that we suppose in the sequel that $j \\neq k$.\n\n\\begin{lemma} \\label{matteo 10}\nFor all $n \\geq 1$, $|l| \\leq N_{n - 1}$, the set $R_{ljk}(i_n) \\subseteq R_{ljk}(i_{n - 1}) $. \n\\end{lemma}\n\n\\begin{proof}\nLike Lemma 5.2 in \\cite{BBM} \n(with $\\omega$ in the role of $\\lambda \\bar\\omega$, and $N_{n-1}$ instead of $N_n$).\n\\end{proof}\n\nBy definition, $ R_{ljk} (i_n) \\subseteq {\\cal G}_n $ (see \\eqref{resonant sets})\nand Lemma \\ref{matteo 10} implies that, for all $ n \\geq 1$, $ |l| \\leq N_{n-1} $, the set\n$R_{ljk}(i_n) \\subseteq R_{ljk}(i_{n - 1}) $. On the other hand \n$ R_{ljk}(i_{n - 1}) \\cap {\\cal G}_{n} = \\emptyset $ (see \\eqref{def:Gn+1}). As a consequence,\nfor all $ |l| \\leq N_{n-1} $, $ R_{ljk} (i_n) = \\emptyset $ and, by \\eqref{Gn inclusioni}, \n\\begin{equation}\\label{Gn Gn+1}\n{\\cal G}_n \\setminus {\\cal G}_{n+1} \\subseteq \\!\\! \\bigcup_{|l| > N_{n-1}, \\, j, k \\in S^c \\cup \\{0\\}} \n\\!\\!\\! \\!\\!\\! R_{ljk} ( i_n) \\qquad \\forall n \\geq 1. \n\\end{equation}\n\n\\begin{lemma}\\label{matteo 4}\nLet $n \\geq 0$. If $R_{ljk}(i_n) \\neq \\emptyset$ \nthen $|l| \\geq C |j^3 - k^3| \\geq \\frac12 C (j^2 + k^2) $ for some $C > 0$.\n\\end{lemma}\n\n\\begin{proof}\nLike Lemma 5.3 in \\cite{BBM}. The only difference is that $\\omega$ is not constrained to a fixed direction. \nNote also that $ | j^3- k^3 | \\geq (j^2 + k^2) \/ 2 $, $ \\forall j \\neq k $. \n\\end{proof}\n\nBy usual arguments (e.g. see Lemma 5.4 in \\cite{BBM}), using Lemma \\ref{matteo 4} and \\eqref{autofinali} we \nhave:\n\n\\begin{lemma} \\label{lemma:risonanti}\nFor all $n \\geq 0$, the measure $ |R_{ljk}(i_n)| \\leq C \\varepsilon^{2(\\nu - 1)} \\gamma \\langle l \\rangle^{- \\tau} $.\n\\end{lemma}\n\n\nBy \\eqref{Gn inclusioni} and Lemmata \\ref{matteo 4}, \\ref{lemma:risonanti} we get\n$$\n|{\\cal G}_0 \\setminus {\\cal G}_1 | \\leq \\sum_{l \\in \\mathbb Z^\\nu, |j|, |k| \\leq C |l|^{1\/2}} \n| R_{ljk} ( i_0)| \\leq \\sum_{l \\in \\mathbb Z^\\nu} \\frac{C \\varepsilon^{2(\\nu-1)} \\gamma }{\\langle l \\rangle^{\\tau-1}} \\leq C' \\varepsilon^{2(\\nu-1)} \\gamma \\,.\n$$\nFor $ n \\geq 1$, by \\eqref{Gn Gn+1}, \n$$\n|{\\cal G}_n \\setminus {\\cal G}_{n+1} | \\leq \\sum_{|l| > N_{n-1}, |j|, |k| \\leq C |l|^{1\/2}} \n| R_{ljk} ( i_n)| \\leq \\sum_{|l| > N_{n-1}} \\frac{C \\varepsilon^{2(\\nu-1)} \\gamma }{\\langle l \\rangle^{\\tau-1}} \\leq C' \\varepsilon^{2(\\nu-1)} \\gamma N_{n-1}^{-1} \n$$\nbecause $\\tau \\geq \\nu + 2$.\nThe estimate $ | \\Omega_\\varepsilon \\setminus {\\cal G}_0| \\leq C \\varepsilon^{2(\\nu-1)} \\gamma $ is elementary. \nThus \\eqref{Gmeasure} is proved. \n\\end{proof}\n\n\n\\noindent\n{\\bf Proof of Theorem \\ref{main theorem} concluded.} \nTheorem \\ref{iterazione-non-lineare} implies that \nthe sequence $ ({\\mathfrak{I}}_n, \\zeta_n) $ \nis well defined for $ \\omega \\in {\\cal G}_\\infty := \\cap_{n \\geq 0} {\\cal G}_n $, that \n$ {\\mathfrak{I}}_n $ \nis a Cauchy sequence in $\\| \\ \\|_{s_0 + \\mu, {\\cal G}_\\infty}^{{\\mathrm{Lip}(\\g)}}$, see \n\\eqref{Hn}, and $ |\\zeta_n |^{\\mathrm{Lip}(\\g)} \\to 0 $. Therefore $ {\\mathfrak{I}}_n $ converges to a limit $ {\\mathfrak{I}}_\\infty $ \nin norm $\\| \\ \\|_{s_0 + \\mu, {\\cal G}_\\infty}^{{\\mathrm{Lip}(\\g)}}$ and, \nby $ ({\\cal P}2)_n $, for all $\\omega \\in {\\cal G}_\\infty$, \n$ i_\\infty (\\varphi) := (\\varphi,0,0) + {\\mathfrak{I}}_\\infty(\\varphi)$, \nis a solution of \n$$\n{\\cal F}(i_\\infty, 0 )= 0\\, \\quad \\text{with} \\quad \n \\| {\\mathfrak{I}}_\\infty \\|^{\\mathrm{Lip}(\\g)}_{s_0 + \\mu, {\\cal G}_\\infty} \n \\leq C \\varepsilon^{6 - 2b} \\gamma^{-1} \n$$\nby \\eqref{ansatz induttivi nell'iterazione} (recall that $b_* := 6 - 2b$). \nTherefore $\\varphi \\mapsto i_\\infty(\\varphi)$ is an invariant torus for the \nHamiltonian vector field $ X_{H_\\varepsilon} $ (see \\eqref{hamiltoniana modificata}).\nBy \\eqref{Gmeasure}, \n$$\n|\\Omega_\\varepsilon \\setminus {\\cal G}_\\infty| \n\\leq |\\Omega_\\varepsilon \\setminus {\\cal G}_0| + \\sum_{n \\geq 0} |{\\cal G}_n \\setminus {\\cal G}_{n + 1}| \n\\leq 2 C_* \\varepsilon^{2(\\nu - 1)} \\gamma + C_* \\varepsilon^{2(\\nu - 1)} \\gamma \\sum_{n \\geq 1} N_{n - 1}^{-1} \n\\leq C \\varepsilon^{2(\\nu - 1)} \\gamma \\,. \n$$\nThe set $\\Omega_\\varepsilon$ in \\eqref{Omega epsilon} has measure $|\\Omega_\\varepsilon | = O( \\varepsilon^{2 \\nu} ) $.\nHence $|\\Omega_\\varepsilon \\setminus \\mathcal{G}_\\infty| \/ |\\Omega_\\varepsilon| \\to 0$ as $\\varepsilon \\to 0$ \nbecause $ \\gamma = o(\\varepsilon^2)$, and therefore the measure of $\\mathcal{C}_\\varepsilon := \\mathcal{G}_{\\infty}$\nsatisfies \\eqref{stima in misura main theorem}.\n\nIn order to complete the proof of Theorem \\ref{main theorem} we show the linear \nstability of the solution $ i_\\infty (\\omega t ) $. \nBy section \\ref{costruzione dell'inverso approssimato} \nthe system obtained linearizing the Hamiltonian vector field $ X_{H_\\varepsilon } $ \nat a quasi-periodic solution $ i_\\infty (\\omega t ) $\nis conjugated to the linear Hamiltonian system\n\\begin{equation}\\label{sistema lineare dopo}\n\\begin{cases}\n\\dot \\psi & \\hspace{-6pt} = K_{20}(\\omega t) \\eta + \nK_{11}^T (\\omega t ) w \n\\\\\n\\dot \\eta & \\hspace{-6pt} = 0 \n\\\\\n\\dot w - \\partial_x K_{0 2}(\\omega t ) w & \\hspace{-6pt} = \\partial_x K_{11}(\\omega t) \\eta \n\\end{cases} \n\\end{equation}\n(recall that the torus $ i_\\infty $ is isotropic and the transformed nonlinear Hamiltonian system is \\eqref{sistema dopo trasformazione inverso approssimato} where $ K_{00}, K_{10}, K_{01} = 0 $, \nsee Remark \\ref{rem:KAM normal form}). \nIn section \\ref{operatore linearizzato sui siti normali} we have proved the reducibility of the linear system\n$ \\dot w - \\partial_x K_{0 2}(\\omega t ) w $, conjugating the last equation in \\eqref{sistema lineare dopo}\nto a diagonal system\n\\begin{equation}\\label{vjmuj}\n{\\dot v}_j + \\mu_j^\\infty v_j = f_j (\\omega t) \\, , \\quad j \\in S^c \\, , \\quad \\mu_j^\\infty \\in {\\mathrm i} \\mathbb R \\, , \n\\end{equation}\nsee \\eqref{Lfinale}, and $ f (\\varphi, x) = \\sum_{j \\in S^c} f_j (\\varphi) e^{{\\mathrm i} j x } \\in H^s_{S^\\bot}\n(\\mathbb T^{\\nu+1}) $. \nThus \\eqref{sistema lineare dopo} is stable.\nIndeed the actions $ \\eta (t) = \\eta_0 \\in \\mathbb R $, $ \\forall t \\in \\mathbb R $. \nMoreover the solutions of the non-homogeneous equation \\eqref{vjmuj} are\n$$\nv_j (t) = c_j e^{ \\mu_j^\\infty t} + {\\tilde v}_j (t) \\,, \\quad \\text{where} \\quad \n{\\tilde v}_j (t) := \\sum_{l \\in \\mathbb Z^\\nu} \\frac{f_{jl} \\, e^{{\\mathrm i} \\omega \\cdot l t } }{ {\\mathrm i} \\omega \\cdot l + \\mu_j^\\infty } \n$$ \nis a quasi-periodic solution (recall that the first Melnikov conditions \\eqref{prime di melnikov} hold at a solution).\nAs a consequence (recall also $ \\mu_j^\\infty \\in {\\mathrm i} \\mathbb R $)\n the Sobolev norm of the solution of \\eqref{vjmuj} with initial condition $ v(0) = \\sum_{j \\in S^c} v_j (0) e^{{\\mathrm i} j x } \n \\in H^{s_0} (\\mathbb T_x) $, $ s_0 < s $, does not increase in time. \\qed\n\\\\[1mm]\n{\\bf Construction of the set $S$ of tangential sites.} \nWe finally prove that, for any $\\nu \\geq 1$, the set $S$ in \\eqref{tang sites} satisfying $({\\mathtt S}1)$-$({\\mathtt S}2)$ can be constructed inductively with only a \\emph{finite} number of restriction at any step of the induction.\n\nFirst, fix any integer $\\bar\\jmath_1 \\geq 1$. \nThen the set $J_1 := \\{ \\pm \\bar\\jmath_1 \\}$ trivially satisfies \n$({\\mathtt S}1) $-$({\\mathtt S}2)$.\nThen, assume that we have fixed $n$ distinct positive integers $\\bar\\jmath_1, \\ldots, \\bar\\jmath_n$, $n \\geq 1$, such that the set $J_n := \\{ \\pm \\bar\\jmath_1, \\ldots, \\pm \\bar\\jmath_n \\}$ satisfies $({\\mathtt S}1)$-$({\\mathtt S}2)$.\nWe describe how to choose another positive integer $\\bar\\jmath_{n+1}$, which is different from all $j \\in J_n$, such that $J_{n+1} := J_n \\cup \\{ \\pm \\bar\\jmath_{n+1} \\}$ also \nsatisfies $({\\mathtt S}1), ({\\mathtt S}2)$.\n\nLet us begin with analyzing $({\\mathtt S}1)$. \nA set of 3 elements $j_1, j_2, j_3 \\in J_{n+1}$ can be of these types: \n$(i)$ all ``old'' elements $j_1, j_2, j_3 \\in J_n$; \n$(ii)$ two ``old'' elements $j_1, j_2 \\in J_n$ and one ``new'' element $j_3 = \\sigma_3 \\bar\\jmath_{n+1}$, $\\sigma_3 = \\pm 1$;\n$(iii)$ one ``old'' element $j_1 \\in J_n$ and two ``new'' elements $j_2 = \\sigma_2 \\bar\\jmath_{n+1}$, $j_3 = \\sigma_3 \\bar\\jmath_{n+1}$, with $\\sigma_2, \\sigma_3 = \\pm 1$;\n$(iv)$ all ``new'' elements $j_i = \\sigma_i \\bar\\jmath_{n+1}$, $\\sigma_i = \\pm 1$, $i = 1,2,3$. \n\n\nIn case $(i)$, the sum $j_1 + j_2 + j_3$ is nonzero by inductive assumption. \nIn case $(ii)$, $j_1 + j_2 + j_3$ is nonzero provided \n$\\bar\\jmath_{n+1} \\notin \\{ j_1 + j_2 : j_1, j_2 \\in J_n \\}$, which is a finite set. \nIn case $(iii)$, for $\\sigma_2 + \\sigma_3 = 0$ the sum $j_1 + j_2 + j_3 = j_1$ is trivially nonzero because $0 \\notin J_n$, while, for $\\sigma_2 + \\sigma_3 \\neq 0$, the sum $j_1 + j_2 + j_3 = j_1 + (\\sigma_2 + \\sigma_3) \\bar\\jmath_{n+1} \\neq 0$ if $\\bar\\jmath_{n+1} \\notin \\{ \\tfrac12 j : j \\in J_n \\}$,\nwhich is a finite set. \nIn case $(iv)$, the sum $j_1 + j_2 + j_3 = (\\sigma_1 + \\sigma_2 + \\sigma_3) \\bar\\jmath_{n+1} \\neq 0$ because \n$\\bar\\jmath_{n+1} \\geq 1$ and $\\sigma_1 + \\sigma_2 + \\sigma_3 \\in \\{ \\pm 1, \\pm 3\\}$. \n\n\\smallskip\n\nNow we study $({\\mathtt S}2)$ for the set $J_{n+1}$. \nDenote, in short, $b := j_1^3 + j_2^3 + j_3^3 + j_4^3 - (j_1 + j_2 + j_3 + j_4)^3$. \n\nA set of 4 elements $j_1, j_2, j_3, j_4 \\in J_{n+1}$ can be of 5 types: \n$(i)$ all ``old'' elements $j_1, j_2, j_3, j_4 \\in J_n$; \n$(ii)$ three ``old'' elements $j_1, j_2, j_3 \\in J_n$ and one ``new'' element $j_4 = \\sigma_4 \\bar\\jmath_{n+1}$, $\\sigma_4 = \\pm 1$;\n$(iii)$ two ``old'' element $j_1, j_2 \\in J_n$ and two ``new'' elements $j_3 = \\sigma_3 \\bar\\jmath_{n+1}$, $j_4 = \\sigma_4 \\bar\\jmath_{n+1}$, with $\\sigma_3, \\sigma_4 = \\pm 1$;\n$(iv)$ one ``old'' element $j_1 \\in J_n$ and three ``new'' elements $j_i = \\sigma_i \\bar\\jmath_{n+1}$, \n$\\sigma_i = \\pm 1$, $i=2,3,4$;\n$(v)$ all ``new'' elements $j_i = \\sigma_i \\bar\\jmath_{n+1}$, $\\sigma_i = \\pm 1$, $i = 1,2,3,4$. \n\nIn case $(i)$, $b \\neq 0$\nby inductive assumption. \n\nIn case $(ii)$, assume that $j_1 + j_2 + j_3 + j_4 \\neq 0$, \nand calculate\n\\begin{align*}\nb & = - 3 (j_1 + j_2 + j_3) \\bar\\jmath_{n+1}^2 \n- 3 (j_1 + j_2 + j_3)^2 \\sigma_4 \\bar\\jmath_{n+1} \n+ [j_1^3 + j_2^3 + j_3^3 - (j_1 + j_2 + j_3)^3] \\ \n=: p_{j_1,j_2,j_3, \\sigma_4}(\\bar\\jmath_{n+1}). \n\\end{align*}\nThis is nonzero provided $p_{j_1,j_2,j_3, \\sigma_4}(\\bar\\jmath_{n+1}) \\neq 0$ for all $j_1, j_2, j_3 \\in J_n$, $\\sigma_4 = \\pm 1$. \nThe polynomial $p_{j_1,j_2,j_3, \\sigma_4}$ is never identically zero because either the leading coefficient $-3(j_1 + j_2 + j_3) \\neq 0$ (and, if one uses $(\\mathtt{S}_3)$, this is always the case), or, if $j_1 + j_2 + j_3 = 0$, then $j_1^3 + j_2^3 + j_3^3 \\neq 0$ \nby \\eqref{prodottino} (using also that $0 \\notin J_n$). \n\n\nIn case $(iii)$, assume that $ j_1 + \\ldots + j_4 = j_1 + j_2 + (\\sigma_3 + \\sigma_4) \\bar\\jmath_{n+1} \\neq 0$, and calculate\n\\begin{align*}\nb & = - 3 \\a \\bar\\jmath_{n+1}^3 \n- 3 \\a^2 (j_1 + j_2) \\bar\\jmath_{n+1}^2\n- 3 (j_1 + j_2)^2 \\a \\bar\\jmath_{n+1} \n- j_1 j_2 (j_1 + j_2) \n=: q_{j_1,j_2,\\a}(\\bar\\jmath_{n+1}), \n\\end{align*}\nwhere $\\a := \\sigma_3 + \\sigma_4$. \nWe impose that $q_{j_1,j_2,\\a}(\\bar\\jmath_{n+1}) \\neq 0$ for all $j_1, j_2 \\in J_n$, \n$\\a \\in\\{ \\pm 2, 0 \\}$. \nThe polynomial $q_{j_1,j_2,\\a}$ is never identically zero because either the leading coefficient $-3\\a \\neq 0$, or, for $\\a = 0$, the constant term $- j_1 j_2 (j_1 + j_2) \\neq 0$ (recall that $0 \\notin J_n$ and $j_1 + j_2 + \\a \\bar\\jmath_{n+1} \\neq 0$). \n\nIn case $(iv)$, assume that $j_1 + \\ldots + j_4 = j_1 + \\a \\bar\\jmath_{n+1} \\neq 0$, \nwhere $\\a := \\sigma_2 + \\sigma_3 + \\sigma_4 \\in \\{ \\pm 1, \\pm 3 \\}$, and calculate\n\\[\nb = \\a \\bar\\jmath_{n+1} r_{j_1,\\a}(\\bar\\jmath_{n+1}), \n\\quad\nr_{j_1, \\a}(x) := (1-\\a^2) x^2 - 3 \\a j_1 x - 3 j_1^2.\n\\]\nThe polynomial $r_{j_1, \\a}$ is never identically zero because $j_1 \\neq 0$. \nWe impose $r_{j_1, \\a}(\\bar\\jmath_{n+1}) \\neq 0$ for all $j_1 \\in J_n$, $\\a \\in \\{ \\pm 1, \\pm 3\\}$. \n\nIn case $(v)$, assume that $j_1 + \\ldots + j_4 = \\a \\bar\\jmath_{n+1} \\neq 0$, with $\\a := \\sigma_1 + \\ldots + \\sigma_4 \\neq 0$, and calculate \n$b = \\a (1 - \\a^2) \\bar\\jmath_{n+1}^3$.\nThis is nonzero because $\\bar\\jmath_{n+1} \\geq 1$ and $\\a \\in \\{\\pm 2, \\pm 4\\}$. \n\nWe have proved that, in choosing $\\bar\\jmath_{n+1}$, there are only finitely many integers to avoid. \n\n\n\n\n\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzoidm b/data_all_eng_slimpj/shuffled/split2/finalzzoidm new file mode 100644 index 0000000000000000000000000000000000000000..ff64ea6d10382fc8683bdd3671017c0763aa46f7 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzoidm @@ -0,0 +1,5 @@ +{"text":"\n\n\n\n\\section{Introduction}\n\\label{Sec:Introduction}\n\nAs numerical simulations of black-hole binaries improve, the criterion\nfor success moves past the ability of a code to merely persist through\nmany orbits of inspiral, merger, and ringdown. Accuracy becomes the\ngoal, as related work in astrophysics and analysis of data from\ngravitational-wave detectors begins to rely more heavily on results\nfrom numerical relativity. One of the most important challenges in\nthe field today is to find and eliminate systematic errors that could\npollute results built on numerics. Though there are many possible\nsources of such error, one stands out as being particularly easy to\nmanage and---as we show---a particularly large effect: the error made\nby extracting gravitational waveforms from a simulation at finite\nradius, and treating these waveforms as though they were the\nasymptotic form of the radiation.\n\nThe desired waveform is the one to which post-Newtonian approximations\naspire, and the one sought by gravitational-wave observatories: the\nasymptotic waveform. This is the waveform as it is at distances of\nover $10^{14}\\,M$ from the system generating the waves. In typical\nnumerical simulations, data extraction takes place at a distance of\norder $100\\,M$ from the black holes. At this radius, the waves are\nstill rapidly changing because of real physical effects. Near-field\neffects~\\cite{Teukolsky1982, Boyle2008, Boyle2009} are plainly\nevident, scaling with powers of the ratio of the reduced wavelength to\nthe radius, $(\\lambdabar\/r)^k$.\\footnote{We use the standard notation\n $\\lambdabar \\equiv \\lambda \/ 2\\pi$.} %\nExtraction methods aiming to eliminate the influence of gauge effects\nalone (\\foreign{e.g.}\\xspace, improved Regge--Wheeler--Zerilli or quasi-Kinnersley\ntechniques) will not be able to account for these physical changes.\n\nEven using a rather naive, gauge-dependent extraction method,\nnear-field effects dominate the error in extracted waves throughout\nthe inspiral for the data presented in this paper~\\cite{Boyle2008}.\nFor extraction at $r=50\\,M$, these effects can account for a\ncumulative error of roughly $50\\%$ in amplitude or a phase difference\nof more than one radian, from beginning to end of a 16-orbit\nequal-mass binary merger. Note that near-field effects should be\nproportional to---at leading order---the ratio of $\\lambdabar\/r$ in\nphase and $(\\lambdabar\/r)^2$ in amplitude, as has been observed\npreviously \\cite{HannamEtAl2008, Boyle2008}. Crucially, because the\nwavelength changes most rapidly during the merger, the amplitude and\nphase differences due to near-field effects also change most rapidly\nduring merger. This means that coherence is lost between the inspiral\nand merger\/ringdown segments of the waveform.\n\nWe can see the importance of this decoherence by looking at its effect\non the matched-filtering technique frequently used to analyze data\nfrom gravitational-wave detectors. Matched filtering~\\cite{Finn1992,\n FinnChernoff1993, BoyleEtAl2009a} compares two signals, $s_{1}(t)$\nand $s_{2}(t)$. It does this by Fourier transforming each into the\nfrequency domain, taking the product of the signals, weighting each\ninversely by the noise---which is a function of frequency---and\nintegrating over all frequencies. This match is optimized over the\ntime and phase offsets of the input waveforms. For appropriately\nnormalized waveforms, the result is a number between 0 and 1, denoted\n$\\Braket{ s_{1} | s_{2}}$, with 0 representing no match, and 1\nrepresenting a perfect match. If we take the extrapolated waveform as\n$s_{1}$ and the waveform extracted at finite radius as $s_{2}$, we can\nevaluate the match between them. If the extrapolated waveform\naccurately represents the ``true'' physical waveform, the mismatch\n(defined as $1-\\Braket{s_{1}|s_{2}}$) shows us the loss of signal in\ndata analysis if we were to use the finite-radius waveforms to search\nfor physical waveforms in detector data.\n\nThe waveforms have a simple scaling with the total mass of the system,\nwhich sets their frequency scale relative to the noise present in the\ndetector. In Figs.~\\ref{fig:MismatchInitial}\nand~\\ref{fig:MismatchAdvanced}, we show mismatches between\nfinite-radius and extrapolated data from the Caltech--Cornell group\nfor a range of masses of interest to LIGO data analysis, using the\nInitial- and Advanced-LIGO noise curves, respectively, to weight the\nmatches. The value of $R$ denotes the coordinate radius of extraction\nfor the finite-radius waveform.\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig1}\n \\caption{\\CapName{Data-analysis mismatch between finite-radius\n waveforms and the extrapolated waveform for Initial LIGO} This\n plot shows the mismatch between extrapolated waveforms and\n waveforms extracted at several finite radii, scaled to various\n values of the total mass of the binary system, using the\n Initial-LIGO noise curve. The waveforms are shifted in time and\n phase to find the optimal match. Note that the data used here is\n solely numerical, with no direct post-Newtonian contribution.\n Thus, for masses below $40\\,\\Sun$, this data represents only a\n portion of the physical waveform.}\n \\label{fig:MismatchInitial}\n\\end{figure}\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig2}\n \\caption{\\CapName{Data-analysis mismatch between finite-radius\n waveforms and the extrapolated waveform for Advanced LIGO} This\n plot shows the mismatch between extrapolated waveforms and\n waveforms extracted at several finite radii, scaled to various\n values of the total mass of the binary system, using the\n Advanced-LIGO noise curve. The waveforms are shifted in time and\n phase to find the optimal match. Note that the data used here is\n solely numerical, with no direct post-Newtonian contribution.\n Thus, for masses below $110\\,\\Sun$, this data represents only a\n portion of the physical waveform.}\n \\label{fig:MismatchAdvanced}\n\\end{figure}\n\nThe data in these figures is exclusively numerical data from the\nsimulation used throughout this paper, with no direct contributions\nfrom post-Newtonian (PN) waveforms. However, to reduce ``turn-on''\nartifacts in the Fourier transforms, we have simply attached PN\nwaveforms to the earliest parts of the time-domain waveforms,\nperformed the Fourier transform, and set to zero all data at\nfrequencies for which PN data are used. The match integrals are\nperformed over the intersection of the frequencies present in each\nwaveform, as in Ref.~\\cite{HannamEtAl2009a}.\n\nThis means that the data used here is not truly complete for masses\nbelow $40\\,\\Sun$ in Initial LIGO and $110\\,\\Sun$ in Advanced LIGO,\nand that a detected signal would actually be dominated by data at\nlower frequencies than are present in this data for masses below about\n$10\\,\\Sun$. These masses are correspondingly larger for shorter\nwaveforms, which begin at higher frequencies. It is important to\nremember that this type of comparison can only show that a given\nwaveform (of a given length) is as good as it needs to be for a\ndetector. If the waveform does not cover the sensitive band of the\ndetector, the detection signal-to-noise ratio would presumably improve\ngiven a comparably accurate waveform of greater duration. Thus, the\nbar is raised for longer waveforms, and for lower masses.\n\nThese figures demonstrate that the mismatch can be of order $1\\%$ when\nextracting at a radius of $R=50\\,M$. For extraction at $R=225\\,M$,\nthe mismatch is never more than about $0.2\\%$. The loss in event rate\nwould be---assuming homogeneous distribution of events in\nspace---roughly 3 times the mismatch when using a template bank based\non imperfect waveforms~\\cite{Brown2004}. Lindblom\n\\foreign{et~al.}\\xspace~\\cite{LindblomEtAl2008} cite a target mismatch of less than\n$0.5\\%$ between the physical waveform and a class of model templates\nto be used for \\emph{detection} of events in current LIGO detector\ndata.\\footnote{ This number of 0.5\\% results from assumptions about\n typical event magnitude, template bank parameters, and requirements\n on the maximum frequency of missed events. The parameters used to\n arrive at this number are typical for Initial LIGO.} Thus, for\nexample, if these numerical waveforms were to be used in construction\nof template banks,\\footnote{ We emphasize that these waveforms do not\n cover the sensitive band of current detectors, and thus would not\n likely be used to construct template banks without the aid of\n post-Newtonian extensions of the data. Longer templates effectively\n have more stringent accuracy requirements, so the suitability of\n these extraction radii would change for waveforms of different\n lengths. In particular, our results are consistent with those of\n Ref.~\\cite{HannamEtAl2009a}, which included non-extrapolated data of\n shorter duration.} the waveform extracted at $R=50\\,M$ would not be\nentirely sufficient, in the sense that a template bank built on\nwaveforms with this level of inaccuracy would lead to an unacceptably\nhigh reduction of event rate. The waveforms extracted at $R=100\\,M$\nand $225\\,M$, on the other hand, may be acceptable for Initial LIGO.\nFor the loudest signals expected to be seen by Advanced LIGO, the\nrequired mismatch may be roughly $10^{-4}$~\\cite{LindblomEtAl2008}.\nIn this case, even extraction at $R=225\\,M$ would be insufficient;\nsome method must be used to obtain the asymptotic waveform. For both\nInitial and Advanced LIGO, estimating the parameters of the\nwaveform---masses and spins of the black holes, for\ninstance---requires still greater accuracy.\n\n\n\nExtrapolation of certain quantities has been used for some time in\nnumerical relativity. Even papers announcing the first successful\nblack-hole binary evolutions \\cite{BuonannoEtAl2007,\n CampanelliEtAl2006, BakerEtAl2006} showed radial extrapolation of\nscalar physical quantities---radiated energy and angular momentum.\nBut waveforms reported in the literature have not always been\nextrapolated. For certain purposes, this is\nacceptable---extrapolation simply removes one of many errors. If the\nprecision required for a given purpose allows it, extrapolation is\nunnecessary. However, for the purposes of LIGO data analysis, we see\nthat extrapolation of the waveform may be very important.\n\nWe can identify three main obstacles to obtaining the asymptotic form\nof gravitational-wave data from numerical simulations:\n\\begin{enumerate}\n \\item Getting the ``right'' data at any given point, independent of\n gauge effects (\\foreign{e.g.}\\xspace, using quasi-Kinnersley techniques and improved\n Regge--Wheeler--Zerilli techniques);\n \\item Removing near-field effects;\n \\item Extracting data along a physically relevant path.\n\\end{enumerate}\nMany groups have attempted to deal with the first of these\nproblems.\\footnote{See~\\cite{NerozziEtAl2005}\n and~\\cite{SarbachTiglio2001} and references therein for descriptions\n of quasi-Kinnersley and RWZ methods, respectively. Also, an\n interesting discussion of RWZ methods, and the possibility of\n finding the ``exact'' waveform at finite distances is found\n in~\\cite{PazosEtAl2007}.} While this is, no doubt, an important\nobjective, even the best extraction technique to date is imperfect at\nfinite radii. Moreover, at finite distances from the source,\ngravitational waves continue to undergo real physical changes as they\nmove away from the system~\\cite{Thorne1980}, which are frequently\nignored in the literature. Some extraction techniques have been\nintroduced that attempt to incorporate corrections for these physical\nnear-field effects~\\cite{AbrahamsEvans1988, AbrahamsEvans1990,\n DeadmanStewart2009}. However, these require assumptions about the\nform of those corrections, which we prefer not to impose. Finally,\neven if we have the optimal data at each point in our spacetime, it is\neasy to see that extraction along an arbitrary (timelike) path through\nthat spacetime could produce a nearly arbitrary waveform, bearing no\nresemblance to a waveform that could be observed in a nearly inertial\ndetector. In particular, if our extraction point is chosen at a\nspecific coordinate location, gauge effects could make that extraction\npoint correspond to a physical path which would not represent any real\ndetector's motion. It is not clear how to estimate the uncertainty\nthis effect would introduce to the waveforms, except by removing the\neffect entirely.\n\nWe propose a simple method using existing data-extraction techniques\nwhich should be able to overcome each of these three obstacles, given\ncertain very basic assumptions. The data are to be extracted at a\nseries of radii---either on a series of concentric spheres, or at\nvarious radii along an outgoing null ray. These data can then be\nexpressed as functions of extraction radius and retarded time using\neither of two simple methods we describe. For each value of retarded\ntime, the waveforms can then be fit to a polynomial in inverse powers\nof the extraction radius. The asymptotic waveform is simply the first\nnonzero term in the polynomial. Though this method also incorporates\ncertain assumptions, they amount to assuming that the data behave as\nradially propagating waves, and that the metric itself is\nasymptotically Minkowski in the coordinates chosen for the simulation.\n\nExtrapolation is, by its very nature, a dangerous procedure. The\nfinal result may be numerically unstable, in the sense that it will\nfail to converge as the order of the extrapolating polynomial is\nincreased. This is to be expected, as the size of the effects to be\nremoved eventually falls below the size of noise in the waveform data.\nThere are likely better methods of determining the asymptotic form of\ngravitational waves produced by numerical simulations. For example,\ncharacteristic evolution is a promising technique that may become\ncommon in the near future~\\cite{BishopEtAl1996, Huebner2001,\n BabiucEtAl2005, BabiucEtAl2008}. Nonetheless, extrapolation does\nprovide a rough and ready technique which can easily be implemented by\nnumerical-relativity groups using existing frameworks.\n\nThis paper presents a simple method for implementing the extrapolation\nof gravitational-wave data from numerical simulations, and the\nmotivation for doing so. In Sec.~\\ref{sec:TortoiseExtrapolation}, we\nbegin by introducing an extrapolation method that uses approximate\ntortoise coordinates, which is the basic method used to extrapolate\ndata in various papers~\\cite{BoyleEtAl2007b, BoyleEtAl2008a,\n ScheelEtAl2009, BoyleEtAl2009a, BuonannoEtAl2009a} by the\nCaltech--Cornell collaboration. The method is tested on the inspiral,\nmerger, and ringdown waveform data of the equal mass, nonspinning,\nquasicircular 15-orbit binary simulation of the Caltech--Cornell\ncollaboration. We present the convergence of the wave phase and\namplitude as the extrapolation order increases, and we also compare\ndata extrapolated using various extraction radii. In\nSec.~\\ref{sec:PhaseExtrapolation}, we propose a different\nextrapolation method using the wave phase---similar to the method\nintroduced in Ref.~\\cite{HannamEtAl2008}---to independently check our\nresults, again demonstrating the convergence properties of the method.\nIn Sec.~\\ref{sec:Comparison}, we compare the extrapolated waveforms of\nboth methods at various extrapolation orders, showing that they agree\nto well within the error estimates of the two methods. A brief\ndiscussion of the pitfalls and future of extrapolation is found in\nSec.~\\ref{sec:Conclusions}. Finally, we include a brief appendix on\ntechniques for filtering noisy data, which is particularly relevant\nhere because extrapolation amplifies noise.\n\n\n\\section{Extrapolation using approximate tortoise coordinates}\n\\label{sec:TortoiseExtrapolation} %\n\nThere are many types of data that can be extracted from a numerical\nsimulation of an isolated source of gravitational waves. The two most\ncommon methods of extracting gravitational waveforms involve using the\nNewman--Penrose $\\Psi_{4}$ quantity, or the metric perturbation $\\ensuremath{h}$\nextracted using Regge--Wheeler--Zerilli techniques. Even if we focus\non a particular type of waveform, the data can be extracted at a\nseries of points along the $z$ axis, for example, or decomposed into\nmultipole components and extracted on a series of spheres around the\nsource. To simplify this introductory discussion of extrapolation, we\nignore the variety of particular types of waveform data. Rather, we\ngeneralize to some abstract quantity $f$, which encapsulates the\nquantity to be extrapolated and behaves roughly as a radially outgoing\nwave.\n\nWe assume that $f$ travels along outgoing null cones, which we\nparametrize by a retarded time $\\ensuremath{t_{\\text{ret}}}$. Along each of these null cones,\nwe further assume that $f$ can be expressed as a convergent (or at\nleast asymptotic) series in $1\/r$---where $r$ is some radial\ncoordinate---for all radii of interest. That is, we assume\n\\begin{equation}\n \\label{eq:FormOfExtrapolatedFunction}\n f(\\ensuremath{t_{\\text{ret}}}, r) = \\sum_{k=0}^{\\infty}\\, \\frac{f_{(k)}(\\ensuremath{t_{\\text{ret}}})} {r^{k}}\\ ,\n\\end{equation}\nfor some functions $f_{(k)}$. The asymptotic behavior of $f$ is given\nby the lowest nonzero $f_{(k)}$.\\footnote{For example, if\n $f=r\\Psi_{4}$, then $f_{(0)}$ gives the asymptotic behavior; if\n $f=\\Psi_{4}$, then $f_{(1)}$ gives the asymptotic behavior.}\n\nGiven data for such an $f$ at a set of retarded times, and a set of\nradii $\\{r_{i}\\}$, it is a simple matter to fit the data for each\nvalue of $\\ensuremath{t_{\\text{ret}}}$ to a polynomial in $1\/r$. That is, for each value of\n$\\ensuremath{t_{\\text{ret}}}$, we take the set of data $\\left\\{ f(\\ensuremath{t_{\\text{ret}}}, r_{i}) \\right\\}$ and\nfit it to a finite polynomial so that\n\\begin{equation}\n \\label{eq:FittingPolynomial}\n f(\\ensuremath{t_{\\text{ret}}}, r_{i}) \\simeq \\sum_{k=0}^{N}\\, \\frac{f_{(k)}(\\ensuremath{t_{\\text{ret}}})}\n {r^{k}_{i}}\\ .\n\\end{equation}\nStandard algorithms~\\cite{PressEtAl2007} can be used to accomplish\nthis fitting; here we use the least-squares method. Of course,\nbecause we are truncating the series of\nEq.~\\eqref{eq:FormOfExtrapolatedFunction} at $k=N$, some of the\neffects from $k>N$ terms will appear at lower orders. We will need to\nchoose $N$ appropriately, checking that the extrapolated quantity has\nconverged sufficiently with respect to this order.\n\n\\subsection{Radial parameter}\n\\label{sec:ChoiceOfR}\nOne subtlety to be considered is the choice of $r$ parameter to be\nused in the extraction and fitting. For numerical simulation of an\nisolated system, one simple and obvious choice is the coordinate\nradius $\\ensuremath{r_{\\text{coord}}}$ used in the simulation. Alternatively, if the data is\nmeasured on some spheroidal surface, it is possible to define an areal\nradius $\\ensuremath{r_{\\text{areal}}}$ by measuring the area of the sphere along with $f$, and\nsetting $\\ensuremath{r_{\\text{areal}}} \\equiv \\sqrt{\\text{area}\/4\\pi}$. Still other choices are\ncertainly possible.\n\nOne objective in choosing a particular $r$ parameter is to ensure the\nphysical relevance of the final extrapolated quantity. If we try to\ndetect the wave, for example, we may want to think of the detector as\nbeing located at some constant value of $r$. Or, we may want $r$ to\nasymptotically represent the luminosity distance. These conditions\nmay be checked by inspecting the asymptotic behavior of the metric\ncomponents in the given coordinates. For example, if the metric\ncomponents in a coordinate system including $r$ asymptotically\napproach those of the standard Minkowski metric, it is not hard to see\nthat an inertial detector could follow a path of constant $r$\nparameter.\n\nSuppose we have two different parameters $r$ and $\\tilde{r}$ which can\nbe related by a series expansion\n\\begin{equation}\n \\label{eq:rRelation}\n r = \\tilde{r}\\, \\left[ 1 + a\/\\tilde{r} + \\ldots \\right]\\ .\n\\end{equation}\nFor the data presented in this paper, we can show that the coordinate\nradius $\\ensuremath{r_{\\text{coord}}}$ and areal radius $\\ensuremath{r_{\\text{areal}}}$ are related in this way.\nIntroducing the expansion coefficients $\\tilde{f}_{(k)}$, we can write\n\\begin{equation}\n \\label{eq:FormOfExtrapolatedFunctionExpanded}\n f(\\ensuremath{t_{\\text{ret}}}, r) = \\sum_{k=0}^{\\infty}\\, \\frac{f_{(k)}(\\ensuremath{t_{\\text{ret}}})} {r^{k}} =\n \\sum_{k=0}^{\\infty}\\, \\frac{\\tilde{f}_{(k)}(\\ensuremath{t_{\\text{ret}}})} {\\tilde{r}^{k}}\\ .\n\\end{equation}\nInserting Eq.~\\eqref{eq:rRelation} into this formula, Taylor\nexpanding, and equating terms of equal order $k$, shows that $f_{(0)}\n= \\tilde{f}_{(0)}$ and $f_{(1)} = \\tilde{f}_{(1)}$. Thus, if the\nasymptotic behavior of $f$ is given by $f_{(0)}$ or $f_{(1)}$, the\nfinal extrapolated data should not depend on whether $r$ or\n$\\tilde{r}$ is used. On the other hand, in practice we truncate these\nseries at finite order. This means that higher-order terms could\n``pollute'' $f_{(0)}$ or $f_{(1)}$. The second objective in choosing\nan $r$ parameter, then, is to ensure fast convergence of the series in\nEq.~\\eqref{eq:FittingPolynomial}. If the extrapolated quantity does\nnot converge quickly as the order of the extrapolating polynomial $N$\nis increased, it may be due to a poor choice of $r$ parameter.\n\nThe coordinate radius used in a simulation may be subject to large\ngauge variations that are physically irrelevant, and hence are not\nreflected in the wave's behavior. That is, the wave may not fall off\nnicely in inverse powers of that coordinate radius. For the data\ndiscussed later in this paper, we find that using the coordinate\nradius of extraction spheres is indeed a poor choice, while using the\nareal radius of those extraction spheres improves the convergence of\nthe extrapolation.\n\n\\subsection{Retarded-time parameter}\n\\label{sec:ChoiceOfRetardedTime} %\n\nSimilar considerations must be made for the choice of retarded-time\nparameter $\\ensuremath{t_{\\text{ret}}}$ to be used in extrapolation. It may be possible to\nevolve null geodesics in numerical simulations, and use these to\ndefine the null curves on which data is to be extracted. While this\nis an interesting possibility that deserves investigation, we propose\ntwo simpler methods here based on an approximate retarded time\nconstructed using the coordinates of the numerical simulation and the\nphase of the waves measured in that coordinate system.\n\nAgain, we have two criteria for choosing a retarded-time parameter.\nFirst is the physical suitability in the asymptotic limit. For\nexample, we might want the asymptotic $\\ensuremath{t_{\\text{ret}}}$ to be (up to an additive\nterm constant in time) the proper time along the path of a detector\nlocated at constant $r$. Again, checking the asymptotic behavior of\nthe metric components with respect to $\\ensuremath{t_{\\text{ret}}}$ and $r$ should be a\nsufficient test of the physical relevance of the parameters. Second,\nwe wish to have rapid convergence of the extrapolation series using\nthe chosen parameter, which also needs to be checked.\n\nAs before, we can also show the equivalence of different choices for\nthe $\\ensuremath{t_{\\text{ret}}}$ parameter. Suppose we have two different approximations\n$\\ensuremath{t_{\\text{ret}}}$ and $\\ensuremath{\\accentset{\\smile}{t}_{\\text{ret}}}$ that can be related by a series expansion\n\\begin{equation}\n \\label{eq:trRelation}\n \\ensuremath{t_{\\text{ret}}} = \\ensuremath{\\accentset{\\smile}{t}_{\\text{ret}}}\\, \\left[ 1 + b\/r + \\ldots \\right]\\ .\n\\end{equation}\nUsing the new expansion coefficients $\\accentset{\\smile}{f}_{(k)}$, we\ncan write\n\\begin{equation}\n \\label{eq:FormOfExtrapolatedFunctionExpandedForTR}\n f(\\ensuremath{t_{\\text{ret}}}, r) = \\sum_{k=0}^{\\infty}\\, \\frac{f_{(k)}(\\ensuremath{t_{\\text{ret}}})} {r^{k}} =\n \\sum_{k=0}^{\\infty}\\, \\frac{\\accentset{\\smile}{f}_{(k)}(\\ensuremath{\\accentset{\\smile}{t}_{\\text{ret}}})}\n {r^{k}}\\ .\n\\end{equation}\nNow, however, we need to assume that the functions $f_{(k)}$ can be\nwell-approximated by Taylor series. If this is true, we can again\nshow that $f_{(0)} = \\accentset{\\smile}{f}_{(0)}$ or, \\emph{if} we\nhave $f_{(0)}=\\accentset{\\smile}{f}_{(0)}=0$, that $f_{(1)} =\n\\accentset{\\smile}{f}_{(1)}$. The condition that $f$ be\nwell-approximated by a Taylor series is nontrivial, and can help to\ninform the choice of $f$. Similarly, the speed of convergence of the\nextrapolation can help to inform the choice of a particular $\\ensuremath{t_{\\text{ret}}}$\nparameter. While it has been shown \\cite{DamourEtAl2008} that a\nretarded-time parameter as simple as $\\ensuremath{t_{\\text{ret}}}=T-R$ is sufficient for some\npurposes, we find that convergence during and after merger is\ndrastically improved when using a somewhat more careful choice.\n\nSince we will be considering radiation from an isolated compact\nsource, our basic model for $\\ensuremath{t_{\\text{ret}}}$ comes from the Schwarzschild\nspacetime; we assume that the system in question approaches this\nspacetime at increasing distance. In analogy with the\ntime-retardation effect on outgoing null rays in a Schwarzschild\nspacetime~\\cite{Chandrasekhar1992}, we define a ``tortoise\ncoordinate'' $\\ensuremath{r_{\\ast}}$ by:\n\\begin{equation}\n \\label{eq:TortoiseCoordinate}\n \\ensuremath{r_{\\ast}} \\equiv r + 2 \\Eadm \\ln\\left(\\frac{r}{2\\Eadm} - 1 \\right)\\\n ,\n\\end{equation}\nwhere $\\Eadm$ is the ADM mass of the initial data.\\footnote{ Kocsis\n and Loeb~\\cite{KocsisLoeb2007} pointed out that the propagation of a\n roughly spherical gravitational wave should be affected primarily by\n the amount of mass \\emph{interior to} the wave. Because the waves\n from a merging binary can carry off a significant fraction\n (typically a few percent) of the binary's mass, this suggests that\n we should allow the mass in this formula to vary in time, falling by\n perhaps a few percent over the duration of the waveform. However,\n this is a small correction of a small correction; we have not found\n it necessary. Perhaps with more refined methods, this additional\n correction would be relevant.} %\nIn standard Schwarzschild coordinates, the appropriate retarded time\nwould be given by $\\ensuremath{t_{\\text{ret}}} = t - \\ensuremath{r_{\\ast}}$. It is not hard to see that the\nexterior derivative $\\d \\ensuremath{t_{\\text{ret}}}$ is null with respect to the Schwarzschild\nmetric.\n\nTaking inspiration from this, we can attempt to account for certain\ndifferences from a Schwarzschild background. Let $T$ and $R$ denote\nthe simulation's coordinates, and suppose that we extract the metric\ncomponents $g^{TT}$, $g^{TR}$, and $g^{RR}$ from the simulation. We\nseek a $\\ensuremath{t_{\\text{ret}}}(T,R)$ such that\n\\begin{equation}\n \\label{eq:dTRetarded}\n \\d\\ensuremath{t_{\\text{ret}}} = \\frac{\\partial \\ensuremath{t_{\\text{ret}}}}{\\partial T}\\, \\d T +\\frac{\\partial\n \\ensuremath{t_{\\text{ret}}}}{\\partial R}\\, \\d R\n\\end{equation}\nis null with respect to these metric components. That is, we seek a\n$\\ensuremath{t_{\\text{ret}}}$ such that\n\\begin{multline}\n \\label{eq:NullCondition}\n g^{TT}\\, \\left( \\frac{\\partial \\ensuremath{t_{\\text{ret}}}}{\\partial T} \\right)^{2} + 2\n g^{TR}\\, \\left( \\frac{\\partial \\ensuremath{t_{\\text{ret}}}}{\\partial T} \\right)\\, \\left(\n \\frac{\\partial \\ensuremath{t_{\\text{ret}}}}{\\partial R} \\right) \\\\ + g^{RR}\\, \\left(\n \\frac{\\partial \\ensuremath{t_{\\text{ret}}}}{\\partial R} \\right)^{2} = 0\\ .\n\\end{multline}\nWe introduce the ansatz $\\ensuremath{t_{\\text{ret}}} = t - \\ensuremath{r_{\\ast}}$, where $t$ is assumed to be a\nslowly varying function of $R$,\\footnote{More specifically, we need\n $\\lvert \\partial t\/\\partial R \\rvert \\ll \\lvert \\partial \\ensuremath{r_{\\ast}}\n \/ \\partial R \\rvert$. This condition needs to be checked for all\n radii used, at all times in the simulation. For the data presented\n below, we have checked this, and shown it to be a valid assumption,\n at the radii used for extrapolation.} %\nand $\\ensuremath{r_{\\ast}}$ is given by Eq.~\\eqref{eq:TortoiseCoordinate} with $R$ in\nplace of $r$ on the right side. If we ignore $\\partial t \/ \\partial\nR$ and insert our ansatz into Eq.~\\eqref{eq:NullCondition}, we have\n\\begin{multline}\n \\label{eq:NullConditionB}\n g^{TT}\\, \\left( \\frac{\\partial t}{\\partial T} \\right)^{2} - 2\n g^{TR}\\, \\left( \\frac{\\partial t}{\\partial T} \\right)\\, \\left(\n \\frac{1}{1-2\\Eadm\/R} \\right) \\\\ + g^{RR}\\, \\left(\n \\frac{1}{1-2\\Eadm\/R} \\right)^{2} = 0\\ .\n\\end{multline}\nWe can solve this for $\\partial t \/ \\partial T$:\n\\begin{equation}\n \\label{eq:RetardedTimeSolutionA}\n \\frac{\\partial t} {\\partial T} = \\frac{1}{1-2\\Eadm\/R}\\,\n \\frac{g^{TR} \\pm \\sqrt{(g^{TR})^{2} - \\, g^{TT}\\, g^{RR}}} {g^{TT}}\n \\ .\n\\end{equation}\nSubstituting the Schwarzschild metric components shows that we should\nchoose the negative sign in the numerator of the second factor.\nFinally, we can integrate (numerically) to find\n\\begin{equation}\n \\label{eq:FullRetardedTimeSolution}\n t = \\int_{0}^{T}\\, \\frac{1}{g^{TT}}\\, \\frac{g^{TR} -\n \\sqrt{(g^{TR})^{2} - g^{TT}\\, g^{RR}}} {1-2\\Eadm\/R}\\, \\d T'\\ .\n\\end{equation}\nNow, in the case where $g^{TR}$ is small compared to 1, we may wish to\nignore it, in which case we have\n\\begin{equation}\n \\label{eq:RetardedTimeSolutionB}\n t = \\int_{0}^{T}\\, \\frac{\\sqrt{-g^{RR} \/ g^{TT}}} {1-2\\Eadm\/R}\\,\n \\d T'\\ .\n\\end{equation}\nIt is not hard to see that this correctly reduces to $t=T$ in the\nSchwarzschild case.\n\nFor the data discussed later in this paper, we make further\nassumptions that $g^{RR} = 1-2\\Eadm\/R$, and that $R=\\ensuremath{r_{\\text{areal}}}$. That is,\nwe define the corrected time\n\\begin{subequations}\n \\label{eq:DynamicLapseCorrection}\n \\begin{gather}\n \\ensuremath{t_{\\text{corr}}} \\equiv \\int_{0}^{T}\\, \\sqrt{\\frac{-1\/g^{TT}}{1 -\n 2\\Eadm\/\\ensuremath{r_{\\text{areal}}}}} \\, \\d T' \\intertext{and the retarded time} \\ensuremath{t_{\\text{ret}}}\n \\equiv \\ensuremath{t_{\\text{corr}}} - \\ensuremath{r_{\\ast}}\\ .\n \\end{gather}\n\\end{subequations}\nWe find that this corrected time leads to a significant improvement\nover the naive choice of $t(T)=T$, while no improvement results from\nusing Eq.~\\eqref{eq:FullRetardedTimeSolution}.\n\n\\subsection{Application to a binary inspiral}\n\\label{sec:Application}\nTo begin the extrapolation procedure, we extract the (spin-weight\n$s=-2$) $(l,m)=(2,2)$ component of $\\Psi_{4}$ data on a set of spheres\nat constant coordinate radius in the simulation.\\footnote{See\n Ref.~\\cite{ScheelEtAl2009} for details of the extraction procedure.\n We use $\\Psi_{4}$ data here, rather than Regge--Wheeler--Zerilli\n data because the $\\Psi_{4}$ data from this simulation is of higher\n quality; it appears that the RWZ data is more sensitive to changes\n in gauge conditions after the merger. This problem is still under\n investigation.} %\nIn the black-hole binary simulations used here (the same as those\ndiscussed in Refs.~\\cite{BoyleEtAl2007b, Boyle2008, BoyleEtAl2008a,\n ScheelEtAl2009}), these spheres are located roughly\\footnote{\n Explicitly, the extraction spheres are at radii $\\ensuremath{r_{\\text{coord}}}\/\\ensuremath{M_{\\text{irr}}} = \\{75,\n 85, 100, 110, 120, \\ldots, 190, 200, 210, 225\\}$, though we find\n that the final result is not sensitive to the exact placement of the\n extraction spheres.} %\nevery $\\Delta\\ensuremath{r_{\\text{coord}}} \\approx 10\\Mirr$ (where $\\Mirr$ is the sum of\nthe irreducible masses of the black holes in the initial data) from an\ninner radius of $\\ensuremath{r_{\\text{coord}}}=75\\Mirr$ to an outer radius of\n$\\ensuremath{r_{\\text{coord}}}=225\\Mirr$, where $\\Mirr$ denotes the total apparent-horizon\nmass of the two holes at the beginning of the simulation. This\nextraction occurs at time steps of $\\Delta \\ensuremath{t_{\\text{coord}}} \\approx 0.5\\Mirr$\nthroughout the simulation. We also measure the areal radius, $\\ensuremath{r_{\\text{areal}}}$,\nof these spheres by integrating the induced area element over the\nsphere to find the area, and defining $\\ensuremath{r_{\\text{areal}}} \\equiv\n\\sqrt{\\text{area}\/4\\pi}$. This typically differs from the coordinate\nradius $\\ensuremath{r_{\\text{coord}}}$ by roughly $\\Mirr\/\\ensuremath{r_{\\text{coord}}}$. Because of gauge effects, the\nareal radius of a coordinate sphere changes as a function of time, so\nwe measure this as a function of time. Finally, we measure the\naverage value of $g^{TT}$ as a function of coordinate time on the\nextraction spheres to correct for the dynamic lapse function. The\nareal radius and $g^{TT}$ are then used to compute the retarded time\n$\\ensuremath{t_{\\text{ret}}}$ defined in Eq.~\\eqref{eq:DynamicLapseCorrection}.\n\nThe gravitational-wave data $\\Psi_{4}$, the areal radius $\\ensuremath{r_{\\text{areal}}}$, and\nthe lapse $N$ are all measured as functions of the code coordinates\n$\\ensuremath{t_{\\text{coord}}}$ and $\\ensuremath{r_{\\text{coord}}}$. We can use these to construct the retarded time\ndefined in Eq.~\\eqref{eq:DynamicLapseCorrection}, using $\\ensuremath{r_{\\text{areal}}}$ in place\nof $r$. This, then, will also be a function of the code coordinates.\nThe mapping between $(\\ensuremath{t_{\\text{ret}}},\\ensuremath{r_{\\text{areal}}})$ and $(\\ensuremath{t_{\\text{coord}}},\\ensuremath{r_{\\text{coord}}})$ is invertible, so we\ncan rewrite $\\Psi_{4}$ as a function of $\\ensuremath{t_{\\text{ret}}}$ and $\\ensuremath{r_{\\text{areal}}}$.\n\nAs noted in Sec.~\\ref{sec:ChoiceOfRetardedTime}, we need to assume\nthat the extrapolated functions are well approximated by Taylor\nseries. Because the real and imaginary parts of $\\Psi_{4}$ are\nrapidly oscillating in the data presented here, we prefer to use the\nsame data in smoother form. We define the complex amplitude $A$ and\nphase $\\phi$ of the wave:\n\\begin{equation}\n \\label{eq:AmplitudeAndPhaseDefinition}\n \n \\ensuremath{r_{\\text{areal}}}\\, \\Mirr\\, \\Psi_4 \\equiv A\\, \\ensuremath{\\mathrm{e}}^{\\i \\phi}\\ ,\n\\end{equation}\nwhere $A$ and $\\phi$ are functions of $\\ensuremath{t_{\\text{ret}}}$ and $\\ensuremath{r_{\\text{areal}}}$. Note that this\ndefinition factors out the dominant $1\/r$ behavior of the amplitude.\nThis equation defines the phase with an ambiguity of multiples of\n$2\\pi$. In practice, we ensure that the phase is continuous as a\nfunction of time by adding suitable multiples of $2\\pi$. The\ncontinuous phase is easier to work with for practical reasons, and is\ncertainly much better approximated by a Taylor series, as required by\nthe argument surrounding\nEq.~\\eqref{eq:FormOfExtrapolatedFunctionExpandedForTR}.\n\nA slight complication arises in the relative phase offset between\nsuccessive radii. Noise in the early parts of the waveform makes the\noverall phase offset go through multiples of $2\\pi$ essentially\nrandomly. We choose some fairly noise-free (retarded) time and ensure\nthat phases corresponding to successive extraction spheres are matched\nat that time, by simply adding multiples of $2\\pi$ to the phase of the\nentire waveform---that is, we add a multiple of $2\\pi$ to the phase at\nall times.\n\nExtrapolation of the waveform, then, basically consists of finding the\nasymptotic forms of these functions, $A$ and $\\phi$ as functions of\ntime. We apply the general technique discussed above to $A$ and\n$\\phi$. Explicitly, we fit the data to polynomials in $1\/\\ensuremath{r_{\\text{areal}}}$ for\neach value of retarded time:\n\\begin{subequations}\n \\label{eq:ExtrapolationFormula}\n \\begin{align}\n \\label{eq:AxmplitudeExtrapolation}\n A(\\ensuremath{t_{\\text{ret}}},\\ensuremath{r_{\\text{areal}}}) &\\simeq \\sum_{k=0}^N\\,\\frac{A_{(k)}(\\ensuremath{t_{\\text{ret}}})}{\\ensuremath{r_{\\text{areal}}}^k}\\ , \\\\\n \\label{eq:PhaseExtrapolation}\n \\phi(\\ensuremath{t_{\\text{ret}}},\\ensuremath{r_{\\text{areal}}}) &\\simeq \\sum_{k=0}^N\\,\\frac{\\phi_{(k)}(\\ensuremath{t_{\\text{ret}}})}{\\ensuremath{r_{\\text{areal}}}^k}\\\n .\n \\end{align}\n\\end{subequations}\nThe asymptotic waveform is fully described by $A_{(0)}$ and\n$\\phi_{(0)}$. When the order of the approximating polynomials is\nimportant, we will denote by $A_{N}$ and $\\phi_{N}$ the asymptotic\nwaveforms resulting from approximations using polynomials of order\n$N$.\n\nWe show the results of these extrapolations in the figures below.\nFigs.~\\ref{fig:LapseCorrectionComparison_Corr_Amp}\nthrough~\\ref{fig:LapseCorrectionComparison_NoCorrMatched} show\nconvergence plots for extrapolations using orders $N=1$--$5$. The\nfirst two figures show the relative amplitude and phase difference\nbetween successive orders of extrapolation, using the corrected time\nof Eq.~\\eqref{eq:DynamicLapseCorrection}. Here, we define\n\\begin{subequations}\n \\begin{gather}\n \\label{eq:RelativeAmplitudeDifferenceDefinition}\n \\frac{\\delta A}{A} \\equiv \\frac{A_{N_{a}} - A_{N_{b}}} {A_{N_{b}}}\n \\\\ \\intertext{and}\n \\label{eq:PhaseDifferenceDefinition}\n \\delta \\phi \\equiv \\phi_{N_{a}} - \\phi_{N_{b}}\\ .\n \\end{gather}\n\\end{subequations}\nWhen comparing waveforms extrapolated by polynomials of different\norders, we use $N_{b}=N_{a}+1$. Note that the broad trend is toward\nconvergence, though high-frequency noise is more evident as the order\nincreases, as we discuss further in the next subsection. The peak\namplitude of the waves occurs at time $\\ensuremath{t_{\\text{ret}}}\/\\ensuremath{M_{\\text{irr}}} \\approx 3954$. Note\nthat the scale of the horizontal axis changes just before this time to\nbetter show the merger\/ringdown portion. We see that the\nextrapolation is no longer convergent, with differences increasing\nslightly as the order of the extrapolating polynomial is increased.\nThe oscillations we see in these convergence plots have a frequency\nequal to the frequency of the waves themselves. Their origin is not\nclear, but may be due to numerics, gauge, or other effects that\nviolate our assumptions about the outgoing-wave nature of the data.\nIt is also possible that there are simply no higher-order effects to\nbe extrapolated, so low-order extrapolation suffices.\n\nFigure~\\ref{fig:LapseCorrectionComparison_NoCorrMatched} shows the\nsame data as in Fig.~\\ref{fig:LapseCorrectionComparison_Corr}, except\nthat no correction is used for dynamic lapse. That is, for this\nfigure (and only this figure), we use $\\ensuremath{t_{\\text{ret}}} \\equiv T - \\ensuremath{r_{\\ast}}$, where $T$\nis simply the coordinate time. This demonstrates the need for\nimproved time-retardation methods after merger. Note that the\nextrapolated data during the long inspiral is virtually unchanged\n(note the different vertical axes). After the merger---occurring at\nroughly $\\ensuremath{t_{\\text{ret}}}\/\\ensuremath{M_{\\text{irr}}} = 3954$---there is no convergence when no\ncorrection is made for dynamic lapse. It is precisely the merger and\nringdown segment during which extreme gauge changes are present in the\ndata used here~\\cite{ScheelEtAl2009}. On the other hand, the fair\nconvergence of the corrected waveforms indicates that it is possible\nto successfully remove these gauge effects.\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig3}\n \\caption{\\CapName{Convergence of the amplitude of the extrapolated\n $\\Psi_{4}$, with increasing order of the extrapolating\n polynomial, $N$} This figure shows the convergence of the\n relative amplitude of the extrapolated Newman--Penrose waveform,\n as the order $N$ of the extrapolating polynomial is increased.\n (See Eq.~\\eqref{eq:ExtrapolationFormula}.) That is, we subtract\n the amplitudes of the two waveforms, and normalize at each time by\n the amplitude of the second waveform. We see that increasing the\n order tends to amplify the apparent noise during the early and\n late parts of the waveform. Nonetheless, the broad\n (low-frequency) trend is towards convergence. Note that the\n differences decrease as the system nears merger; this is a first\n indication that the extrapolated effects are due to near-field\n influences. Also note that the horizontal axis changes in the\n right part of the figure, which shows the point of merger, and the\n ringdown portion of the waveform. After the merger, the\n extrapolation is nonconvergent, though the differences grow slowly\n with the order of extrapolation.}\n \\label{fig:LapseCorrectionComparison_Corr_Amp}\n\\end{figure}\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig4}\n \\caption{\\CapName{Convergence of the phase of the extrapolated\n $\\Psi_{4}$, with increasing order of the extrapolating\n polynomial, $N$} This figure is the same as\n Fig.~\\ref{fig:LapseCorrectionComparison_Corr_Amp}, except that it\n shows the convergence of phase. Again, increasing the\n extrapolation order tends to amplify the noise during the early\n and late parts of the waveform, though the broad (low-frequency)\n trend is towards convergence. The horizontal-axis scale changes\n just before merger.}\n \\label{fig:LapseCorrectionComparison_Corr}\n\\end{figure}\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig5}\n \\caption{\\CapName{Convergence of the phase of $\\Psi_{4}$,\n extrapolated with no correction for the dynamic lapse} This\n figure is the same as\n Fig.~\\ref{fig:LapseCorrectionComparison_Corr}, except that no\n correction is made to account for the dynamic lapse. (See\n Eq.~\\eqref{eq:DynamicLapseCorrection} and surrounding discussion.)\n Observe that the convergence is very poor after merger (at roughly\n $\\ensuremath{t_{\\text{ret}}}\/\\ensuremath{M_{\\text{irr}}} = 3954$). This corresponds to the time after which\n sharp features in the lapse are observed. We conclude from this\n graph and comparison with the previous graph that the correction\n is crucial to convergence of $\\Psi_{4}$ extrapolation through\n merger and ringdown.}\n \\label{fig:LapseCorrectionComparison_NoCorrMatched}\n\\end{figure}\n\n\n\\subsection{Choosing the order of extrapolation}\n\\label{sec:ExtrapolationOrder} %\nDeciding on an appropriate order of extrapolation to be used for a\ngiven purpose requires balancing competing effects. As we see in\nFig.~\\ref{fig:LapseCorrectionComparison_Corr_Amp}, for example, there\nis evidently some benefit to be gained from using higher-order\nextrapolation during the inspiral; there is clearly some convergence\nduring inspiral for each of the orders shown. On the other hand,\nhigher-order methods amplify the apparent noise in the\nwaveform.\\footnote{So-called ``junk radiation'' is a ubiquitous\n feature of initial data for current numerical simulations of binary\n black-hole systems. It is clearly evident in simulations as\n large-amplitude, high-frequency waves that die out as the simulation\n progresses. While it is astrophysically extraneous, it is\n nevertheless a correct result of evolution from the initial data.\n Better initial data would, presumably, decrease its magnitude. This\n is the source of what looks like noise in the waveforms at early\n times. It is less apparent in $\\ensuremath{h}$ data than in $\\Psi_{4}$ data\n because $\\Psi_{4}$ effectively amplifies high-frequency components,\n because of the relation $\\Psi_{4} \\approx -\\ddot{h}$.} %\nMoreover, late in the inspiral, and on into the merger and ringdown,\nthe effects being extrapolated may be present only at low orders;\nincreasing the extrapolation order would be useless as higher-order\nterms would simply be fitting to noise.\n\nThe optimal order depends on the accuracy needed, and on the size of\neffects that need to be eliminated from the data. For some\napplications, little accuracy is needed, so a low-order extrapolation\n(or even no extrapolation) is preferable.\\footnote{ We note that---as\n expected from investigations of near-field effects\n \\cite{Teukolsky1982, Boyle2008, Boyle2009}---the second-order\n behavior of the amplitude greatly dominates its first-order behavior\n \\cite{HannamEtAl2008}. Thus, there is no improvement to the\n accuracy \\emph{of the amplitude} when extrapolating with $N=1$; it\n would be better to simply use the data from the largest extraction\n radius.} If high-frequency noise is not considered a problem, then\nsimple high-order extrapolation should suffice. Of course, if both\nhigh accuracy and low noise are required, data may easily be filtered,\nmitigating the problem of noise amplification. (See the appendix for\nmore discussion.) There is some concern that this may introduce\nsubtle inaccuracies: filtering is more art than science, and it is\ndifficult to establish precise error bars for filtered data.\n\n\n\\subsection{Choosing extraction radii}\n\\label{sec:ExtractionRadii} %\n\nAnother decision needs to be made regarding the number and location of\nextraction surfaces. Choosing the number of surfaces is fairly easy,\nbecause there is typically little cost in increasing the number of\nextraction radii (especially relative to the cost of---say---running a\nsimulation). The only restriction is that the number of data points\nneeds to be significantly larger than the order of the extrapolating\npolynomial; more can hardly hurt. More careful consideration needs to\nbe given to the \\emph{location} of the extraction surfaces.\n\nFor the extrapolations shown in\nFigs.~\\ref{fig:LapseCorrectionComparison_Corr_Amp}\nand~\\ref{fig:LapseCorrectionComparison_Corr}, data was extracted on\nspheres spaced by $10$ to $15\\Mirr$, from $R=75\\Mirr$ to\n$R=225\\Mirr$. The outer radius of $225\\Mirr$ was chosen simply\nbecause this is the largest radius at which data exists throughout the\nsimulation; presumably, we always want the outermost radii at which\nthe data are resolved. In choosing the inner radius, there are two\ncompeting considerations.\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig6}\n \\caption{\\CapName{Comparison of extrapolation of $\\Psi_{4}$ using\n different sets of extraction radii} This figure compares the\n phase of waveforms extrapolated with various sets of radii. All\n comparisons are with respect to the data set used elsewhere in\n this paper, which uses extraction radii $R\/\\Mirr = \\{75, 85,\n 100, 110, 120, \\ldots, 200, 210, 225\\}$. The order of the\n extrapolating polynomial is $N=3$ in all cases.}\n \\label{fig:RadiiComparison}\n\\end{figure}\n\nOn one hand, we want the largest spread possible between the inner and\nouter extraction radii to stabilize the extrapolation. A very rough\nrule of thumb seems to be that the distance to be extrapolated should\nbe no greater than the distance covered by the data. Because the\nextrapolating polynomial is a function of $1\/R$, the distance to be\nextrapolated is $1\/R_{\\text{outer}} - 1\/\\infty = 1\/R_{\\text{outer}}$.\nThe distance covered by the data is $1\/R_{\\text{inner}} -\n1\/R_{\\text{outer}}$, so if the rule of thumb is to be satisfied, the\ninner extraction radius should be no more than half of the outer\nextraction radius, $R_{\\text{inner}} \\lesssim R_{\\text{outer}}\/2$\n(noting, of course, that this is a \\emph{very} rough rule of thumb).\n\nOn the other hand, we would like the inner extraction radius to be as\nfar out as possible. Extracting data near the violent center of the\nsimulation is a bad idea for many reasons. Coordinate ambiguity,\ntetrad errors, near-field effects---all are more severe near the\ncenter of the simulation. The larger these errors are, the more work\nthe extrapolation needs to do. This effectively means that\nhigher-order extrapolation is needed if data are extracted at small\nradii. The exact inner radius needed for extrapolation depends on the\ndesired accuracy and, again, the portion of the simulation from which\nthe waveform is needed.\n\nWe can compare data extrapolated using different sets of radii.\nFigure~\\ref{fig:RadiiComparison} shows a variety, compared to the data\nused elsewhere in this paper. The extrapolation order is $N=3$ in all\ncases. Note that the waveforms labeled $R\/\\Mirr = \\{50, \\ldots,\n100\\}$ and $R\/\\Mirr = \\{100, \\ldots, 225\\}$ both satisfy the rule\nof thumb that the inner radius should be at most half of the outer\nradius, while the other two waveforms do not; it appears that\nviolation of the rule of thumb leads to greater sensitivity to noise.\nOne waveform is extrapolated using only data from small radii,\n$R\/\\Mirr = \\{50, \\ldots, 100\\}$. It is clearly not converged, and\nwould require higher-order extrapolation if greater accuracy is\nneeded. The source of the difference is presumably the near-field\neffect~\\cite{Boyle2008}, which is proportionally larger at small\nradii.\n\nClearly, there is a nontrivial interplay between the radii used for\nextraction and the order of extrapolation. Indeed, because of the\ntime-dependence of the various elements of these choices, it may be\nadvisable to use different radii and orders of extrapolation for\ndifferent time portions of the waveform. The different portions could\nthen be joined together using any of various\nmethods~\\cite{AjithEtAl2008, BoyleEtAl2009a}.\n\n\n\\section{Extrapolation using the phase of the waveform}\n\\label{sec:PhaseExtrapolation}\n\nWhile the tortoise-coordinate method just described attempts to\ncompensate for nontrivial gauge perturbations, it is possible that it\ndoes not take account of all effects adequately. As an independent\ncheck, we discuss what is essentially a second---very\ndifferent---formulation of the retarded-time parameter, similar to one\nfirst introduced in Ref.~\\cite{HannamEtAl2008}. If waves extrapolated\nwith the two different methods agree, then we can be reasonably\nconfident that unmodeled gauge effects are not diminishing the\naccuracy of the final result. As we will explain below, the method in\nthis section cannot be used naively with general data (\\foreign{e.g.}\\xspace, data on\nthe equatorial plane). In particular, we must assume that the data to\nbe extrapolated consists of a strictly monotonic phase. It is,\nhowever, frequently possible to employ a simple technique to make\npurely real, oscillating data into complex data with strictly\nmonotonic phase, as we describe below. The results of this technique\nagree with those of the tortoise-coordinate extrapolation as we show\nin Sec.~\\ref{sec:Comparison}.\n\nInstead of extrapolating the wave phase $\\phi$ and amplitude $A$ as\nfunctions of time and radius, we extrapolate the time $\\ensuremath{t_{\\text{ret}}}$ and the\namplitude $A$ as functions of wave phase $\\phi$ and radius $\\ensuremath{r_{\\text{areal}}}$. In\nother words, we measure the amplitude and the arrival time to some\nradius $\\ensuremath{r_{\\text{areal}}}$ of a fixed phase point in the waveform. This is the\norigin of the requirement that the data to be extrapolated consist of\na strictly monotonic phase $\\phi(\\ensuremath{t_{\\text{ret}}}, \\ensuremath{r_{\\text{areal}}})$ (\\foreign{i.e.}\\xspace, it must be\ninvertible). For the data presented here, the presence of radiation\nin the initial data---junk radiation---and numerical noise cause the\nextracted waveforms to fail to satisfy this requirement at early\ntimes. In this case, the extrapolation is performed separately for\neach invertible portion of the data. That is, the data are divided\ninto invertible segments, each segment is extrapolated separately, and\nthe final products are joined together as a single waveform.\n\n\n\\subsection{Description of the method}\n\\label{sec:PhaseDescription} %\nThis extrapolation technique consists of extrapolating the retarded\ntime and the amplitude as functions of the wave phase $\\phi$ and the\nradius $\\ensuremath{r_{\\text{areal}}}$. In other words, when extrapolating the waveform, we are\nestimating the amplitude and the arrival time of a fixed phase point\nat infinity. Here, we extract the same $\\Psi_{4}$, $g^{TT}$, and\nareal-radius data used in the previous section. As in the previous\nmethod, we first shift each waveform in time using $\\ensuremath{t_{\\text{ret}}} = \\ensuremath{t_{\\text{corr}}} -\n\\ensuremath{r_{\\ast}}$, where $\\ensuremath{t_{\\text{corr}}}$ is defined in\nEq.~\\eqref{eq:DynamicLapseCorrection} and the basic tortoise\ncoordinate $\\ensuremath{r_{\\ast}}$ is defined in Eq.~\\eqref{eq:TortoiseCoordinate} with\nareal radius as the radial parameter. The amplitude and wave phase\nare again defined using Eq.~\\eqref{eq:AmplitudeAndPhaseDefinition},\nand the phase is made continuous as in Sec.~\\ref{sec:Application}.\nThus, we begin with the same data, shifted as with the\ntortoise-coordinate method.\n\nNow, however, we change the method, in an attempt to allow for\nunmodeled effects. Instead of extrapolating $\\phi(\\ensuremath{t_{\\text{ret}}}, \\ensuremath{r_{\\text{areal}}})$ and\n$A(\\ensuremath{t_{\\text{ret}}}, \\ensuremath{r_{\\text{areal}}})$, as with the previous method, we invert these functions\nto get $\\ensuremath{t_{\\text{ret}}}(\\phi,\\ensuremath{r_{\\text{areal}}})$ and $A(\\phi,\\ensuremath{r_{\\text{areal}}})$ as functions of the wave\nphase $\\phi$. In other words, we extrapolate the arrival time and the\namplitude of a signal to a coordinate radius $R$ for each wave phase\nvalue. This is done by fitting the retarded time $\\ensuremath{t_{\\text{ret}}}$ and the\namplitude $A$ data to polynomials in $1\/\\ensuremath{r_{\\text{areal}}}$ for each value of the\nwave phase:\n\\begin{subequations}\n \\label{eq:ExtrapolationFormulaMethod2}\n \\begin{align}\n \\label{eq:AmplitudeExtrapolationMethod2}\n A(\\ensuremath{r_{\\text{areal}}}, \\phi) &\\simeq \\sum^N_{k=0}\\frac{A_{(k)}(\\phi)}{\\ensuremath{r_{\\text{areal}}}^{k}}\n \\ , \\\\\n \\label{eq:PhaseExtrapolationMethod2}\n t(\\ensuremath{r_{\\text{areal}}}, \\phi) &\\simeq \\ensuremath{r_{\\ast}} + \\sum^N_{k=0}\\frac{t_{(k)}(\\phi)}{\\ensuremath{r_{\\text{areal}}}^k}\n \\ ,\n \\end{align}\n\\end{subequations}\nwhere the asymptotic waveform is fully described by $A_{(0)}(\\phi)$\nand $t_{(0)}(\\phi)$.\n\nWith this data in hand, we can produce the asymptotic amplitude and\nphase as functions of time by plotting curves in the \\mbox{$t$--$A$}\nand \\mbox{$t$--$\\phi$} planes parametrized by the phase. In order to\nbe true, single-valued functions, we again need monotonicity of the\n$t_{(0)}(\\phi)$ data, which may be violated by extrapolation. The\nusable data can be obtained simply by removing data from times before\nwhich this condition holds.\n\nChoosing the extraction radii and extrapolation order for this method\nfollows the same rough recommendations described in\nSecs.~\\ref{sec:ExtrapolationOrder} and~\\ref{sec:ExtractionRadii}.\nNote also that the restriction that the data have an invertible phase\nas a function of time is not insurmountable. For example, data for\n$\\Psi_{4}$ in the equatorial plane is purely real, hence has a phase\nthat simply jumps from $0$ to $\\pi$ discontinuously. However, we can\ndefine a new quantity\n\\begin{equation}\n w(t) \\equiv \\Psi_{4}(t) + \\i \\dot{\\Psi}_{4}(t)\\ .\n\\end{equation}\nThis is simply an auxiliary quantity used for the extrapolation, with\na smoothly varying, invertible phase. The imaginary part is discarded\nafter extrapolation.\n\n\n\\subsection{Results}\n\\label{sec:PhaseResults} %\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig7}\n \\caption{\\CapName{Convergence of the amplitude of $\\Psi_{4}$\n extrapolated using the wave phase, with increasing order $N$ of\n the extrapolating polynomial} This figure shows the convergence\n of the relative amplitude of the extrapolated Newman--Penrose\n waveform extrapolated using the wave phase, as the order $N$ of\n the extrapolating polynomial is increased. (See\n Eq.~\\eqref{eq:ExtrapolationFormulaMethod2}.) Increasing the\n extrapolation order tends to amplify the apparent noise during the\n early and late parts of the waveform, but it improves convergence.\n The vertical axis at $\\ensuremath{t_{\\text{ret}}}\/\\ensuremath{M_{\\text{irr}}}\\approx 3950$ denotes merger.}\n \\label{fig:AmplitudeComparisonMethod2}\n\\end{figure}\n\nIn Figs.~\\ref{fig:AmplitudeComparisonMethod2}\nand~\\ref{fig:PhaseComparisonMethod2} we plot the convergence of the\nrelative amplitude and phase of the extrapolated $(l,m) = (2,2)$ mode\nof the $\\Psi_4$ waveform for extrapolation orders $N=1, \\ldots, 5$. A\ncommon feature of both plots is that during the inspiral, the higher\nthe extrapolation order, the better the convergence. However, the\nnoise is amplified significantly for large orders of extrapolation.\nThis method of extrapolation amplifies high-frequency noise\nsignificantly, compared to the tortoise-coordinate method.\n\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig8}\n \\caption{\\CapName{Convergence of the phase of $\\Psi_{4}$ as a\n function of time extrapolated using the wave phase, with\n increasing order $N$ of the extrapolating polynomial} Again,\n increasing the extrapolation order tends to amplify the apparent\n noise during the early and late parts of the waveform, though\n convergence is improved significantly.}\n \\label{fig:PhaseComparisonMethod2}\n\\end{figure}\n\nIn the inspiral portion, we have a decreasing error in the\nextrapolation of the phase and the amplitude as the wavelength of the\ngravitational waves decreases. In the merger\/ringdown portion, a more\ncareful choice of the radii and order of extrapolation needs to be\nmade. Since near-field effects are less significant in the data\nextracted at larger radii, extrapolation at low order ($N=2,3$) seems\nsufficient. Data extrapolated at large order ($N=4,5$) has a larger\nerror in the phase and amplitude after merger than data extrapolated\nat order $N=2$ or $3$. Moreover, the outermost extraction radius\ncould be reduced, say to $R_{\\text{outer}} \/ \\ensuremath{M_{\\text{irr}}} = 165$ instead of\n$R_{\\text{outer}} \/ \\ensuremath{M_{\\text{irr}}} = 225$, without having large extrapolation\nerror at late times. Using the radius range $R\/\\ensuremath{M_{\\text{irr}}}={75,\\ldots,160}$\ninstead of the range $R\/ \\ensuremath{M_{\\text{irr}}} = {75,\\ldots,225}$ would leave the\nextrapolation error during the merger\/ringdown almost unchanged, while\nthe extrapolation error during the inspiral would increase by about\n70\\%.\n\nWe note that this method allows easy extrapolation of various portions\nof the waveform using different extraction radii and orders since---by\nconstruction---the wave phase is an independent variable. For example,\nsolve for the phase value of the merger $\\phi_{\\text{merger}}$\n(defined as the phase at which the amplitude is a maximum), then use\nthe radius range $R \/ \\ensuremath{M_{\\text{irr}}} = {75,\\ldots,225}$ for all phase values\nless than $\\phi_{\\text{merger}}$ and the range\n$R\/\\ensuremath{M_{\\text{irr}}}={75,\\ldots,160}$ for all larger phase values.\n\nThis method has been tested also using the coordinate radius $R$ and\nthe naive time coordinate $T$, in place of areal radius and corrected\ntime. We found results similar to those discussed in\nSec.~\\ref{sec:TortoiseExtrapolation}. Using the new time coordinate\n$\\ensuremath{t_{\\text{corr}}}$ instead of the naive time coordinate $T$ improved the\nextrapolation during the merger\/ringdown, as found in\nSec.~\\ref{sec:TortoiseExtrapolation}.\n\nAs with the previous extrapolation method, increasing the\nextrapolation order gives a faster convergence rate of waveform phase\nand amplitude, but it amplifies noise in the extrapolated waveform. To\nimprove convergence without increasing the noise, we need a good\nfiltering technique for the inspiral data. The junk-radiation noise\ndecreases significantly as a function of time, disappearing several\norbits before merger. However, this noise could be reduced by using\nmore extraction radii in the extrapolation process, or by running the\ndata through a low-pass filter. See the appendix for further\ndiscussion of filtering.\n\n\n\n\\section{Comparison of the two methods}\n\\label{sec:Comparison}\n\nBoth methods showed good convergence of the amplitude and the phase of\nthe waveform as $N$ increased in the inspiral portion. (See\nFigs.~\\ref{fig:LapseCorrectionComparison_Corr_Amp}\nand~\\ref{fig:AmplitudeComparisonMethod2} for the amplitude, and\nFigs.~\\ref{fig:LapseCorrectionComparison_Corr}\nand~\\ref{fig:PhaseComparisonMethod2} for the phase.) The wave-phase\nextrapolation method was more sensitive to noise. In the\nmerger\/ringdown portion, both methods have similar convergence as $N$\nincreases, especially when the correction is taken to account for the\ndynamic lapse. The use of the time parameter $\\ensuremath{t_{\\text{corr}}}$ improved the\nagreement between the methods significantly in the merger\/ringdown\nportion for all extrapolation orders. Extrapolating at order $N=2$ or\n$3$ seems the best choice as the noise and phase differences are\nsmallest for these values.\n\nIn Fig.~\\ref{fig:MroueBoyleRelativeAmplitudeDifferenceByOrder}, we\nshow the relative amplitude difference between data extrapolated at\nvarious orders ($N=1,\\ldots,5$). There is no additional time or phase\noffset used in the comparison. Ignoring high-frequency components,\nthe difference in the relative amplitude is always less than 0.3\\% for\ndifferent extrapolation orders. Even including high-frequency\ncomponents, the differences between the two methods are always smaller\nthan the error in each method, as judged by convergence plots. In\nFig.~\\ref{fig:MroueBoylePhaseDifferenceByOrder}, we show the\n\\emph{phase} difference between the data extrapolated using both\nmethods. As in the relative amplitude-difference plots, the best\nagreement is achieved during the inspiral portion. Ignoring\nhigh-frequency components, the difference is less than 0.02 radians\nfor all orders. In the merger\/ringdown portion, the best agreement\nbetween extrapolated waveforms is at order $N=2$ or $3$ where the\nphase difference is less than 0.01 radians.\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig9}\n \\caption{\\CapName{Relative difference in the amplitude of the two\n extrapolation methods as we increase the order of extrapolation}\n The best agreement between both methods is at low orders of\n extrapolation, for which the relative difference in the amplitude\n is less than 0.1\\% during most of the evolution.}\n \\label{fig:MroueBoyleRelativeAmplitudeDifferenceByOrder}\n\\end{figure}\n\n\\begin{figure}\n \n \\includegraphics[width=\\linewidth]{Fig10}\n \\caption{\\CapName{Phase difference of the two extrapolation methods\n as we increase the order of extrapolation} This figure shows the\n phase difference between waveforms extrapolated using each of the\n two methods. The best agreement between the methods is at orders\n $N=2$ and $3$. The relative difference in the phase is less than\n 0.02 radians during most of the evolution.}\n \\label{fig:MroueBoylePhaseDifferenceByOrder}\n\\end{figure}\n\n\n\n\n\n\\section{Conclusions}\n\\label{sec:Conclusions}\n\nWe have demonstrated two simple techniques for extrapolating\ngravitational-wave data from numerical-relativity simulations. We\ntook certain basic gauge information into account to improve\nconvergence of the extrapolation during times of particularly dynamic\ngauge, and showed that the two methods agree to within rough error\nestimates. We have determined that the first method presented here is\nless sensitive to noise, and more immediately applies to arbitrary\nwavelike data; this method has become the basic standard in use by the\nCaltech--Cornell collaboration. In both cases, there were problems\nwith convergence after merger. The source of these problems is still\nunclear, but will be a topic for further investigation.\n\nAs with any type of extrapolation, a note of caution is in order. It\nis entirely possible that the ``true'' function being extrapolated\nbears little resemblance to the approximating function we choose,\noutside of the domain on which we have data. We may, however, have\nreason to believe that the true function takes a certain form. If the\ndata in question are generated by a homogeneous wave equation, for\ninstance, we know that well-behaved solutions fall off in powers of\n$1\/r$. In any case, there is a certain element of faith that\nextrapolation is a reasonable thing to do. While that faith may be\nmisplaced, there are methods of checking whether or not it is:\ngoodness-of-fit statistics, error estimates, and convergence tests.\nTo be of greatest use, goodness-of-fit statistics and error estimates\nfor the output waveform require error estimates for the input\nwaveforms. We leave this for future work.\n\nWe still do not know the correct answers to the questions numerical\nrelativity considers. We have no analytic solutions to deliver the\nwaveform that Einstein's equations---solved perfectly---predict would\ncome from a black-hole binary merger; or the precise amount of energy\nradiated from any given binary; or the exact kick or spin of the final\nblack hole. Without being able to compare numerical relativity to\nexact solutions, we may be leaving large systematic errors hidden in\nplain view. To eliminate them, we need to use multiple, independent\nmethods for our calculations. For example, we might extract $\\Psi_4$\ndirectly by calculating the Riemann tensor and contracting\nappropriately with our naive coordinate tetrad, and extract the metric\nperturbation using the formalism of Regge--Wheeler--Zerilli and\nMoncrief. By differentiating the latter result twice and comparing to\n$\\Psi_4$, we could verify that details of the extraction methods are\nnot producing systematic errors. (Just such a comparison was done in\nRef.~\\cite{BuonannoEtAl2009a} for waveforms extrapolated using the\ntechnique in this paper.) Nonetheless, it is possible that\ninfrastructure used to find both could be leading to errors.\n\nIn the same way, simulations need to be performed using different\ngauge conditions, numerical techniques, code infrastructures, boundary\nconditions, and even different extrapolation methods. Only when\nmultiple schemes arrive at the same result can we be truly confident\nin any of them. But to arrive at the same result, the waveforms from\neach scheme need to be processed as carefully as possible. We have\nshown that extrapolation is crucial for highly accurate gravitational\nwaveforms, and for optimized detection of mergers in detector data.\n\n\n\n\\begin{acknowledgments}\n We thank Emanuele Berti, Duncan Brown, Luisa Buchman, Alessandra\n Buonanno, Yanbei Chen, \\'{E}anna Flanagan, Mark Hannam, Sascha Husa,\n Luis Lehner, Geoffrey Lovelace, Andrea Nerozzi, Rob Owen, Larne\n Pekowsky, Harald Pfeiffer, Oliver Rinne, Uli Sperhake, B\\'{e}la\n Szil\\'{a}gyi, Kip Thorne, Manuel Tiglio, and Alan Weinstein for\n helpful discussions. We especially thank Larry Kidder, Lee\n Lindblom, Mark Scheel, and Saul Teukolsky for careful readings of\n this paper in various draft forms and helpful comments. This\n research has been supported in part by a grant from the Sherman\n Fairchild Foundation to Caltech and Cornell; by NSF Grants\n PHY-0652952, PHY-0652929 and DMS-0553677 and NASA Grant NNX09AF96G\n to Cornell; and by NSF grants PHY-0601459, PHY-0652995, DMS-0553302\n and NASA grant NNG05GG52G to Caltech.\n\\end{acknowledgments}\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{introduction}\nWhen we perform a quantum measurement and extract information from a system, the system is disturbed due to the backaction of the measurement.\nThe more information we extract by quantum measurement, the more strongly the system is disturbed.\nSuch a tradeoff relation has been recognized since Heisenberg pointed out the uncertainty relation~\\cite{He27} between error and disturbance.\nSince 1990's, with the tremendous development of the quantum information theory, various methods of quantifying the information and the disturbance has been proposed, and the tradeoff relations has been shown based on their definitions\\cite{Fu96,Da03,Sa06,BS06,BHH08,CL12}.\nIn most of the studies, the information and the disturbance are defined in the information theoretic settings, i.e., decoding the message encoded in quantum states by a measurement.\n\nIn this paper, we formulate the information and the disturbance in the setting of estimating an unknown quantum state by quantum measurements, with an emphasis on the estimation accuracy.\nWe derive an inequality that shows the tradeoff relation between them, and give a sufficient condition to achieve the equality.\nWe also discuss the condition for divergences, which measure the distinguishability between two quantum states, to satisfy a similar inequality.\n\n\\section{information-disturbance relation based on estimation theory}\nSuppose that we estimate an unknown quantum state corresponding to a density operator $\\h\\rho_{\\bm\\theta}$ by performing a quantum measurement.\nHere, $\\bm\\theta\\in\\Theta\\subset\\mathbb R^m$ represents $m$ real parameters that characterize the unknown state, so that estimating the state is equivalent to estimating the parameters.\nSuch a parameterized family of states $\\{\\h\\rho_{\\bm\\theta}\\}_{\\bm\\theta\\in\\Theta}$ is called a quantum statistical model.\nA quantum measurement is characterized by a mapping from quantum states to the probability distribution of outcomes and the post-measurement state corresponding to each outcome.\nIf the measurement outcome is discrete, the probability $p_{\\bm\\theta,i}$ of obtaining the outcome $i\\in I$ and the post-measurement state $\\h\\rho_{\\bm\\theta,i}$ are respectively given by \n\\e{\np_{\\bm\\theta,i}&=\\sum_j \\tr{\\h K_{ij}\\h\\rho_{\\bm\\theta}\\h K_{ij}^\\dagger}\\\\\n\\h\\rho_{\\bm\\theta,i}&=\\frac{1}{p_{\\bm\\theta,i}}\\sum_j \\h K_{ij}\\h\\rho_{\\bm\\theta}\\h K_{ij}^\\dagger,\n}\nwhere measurement operators $\\{\\h K_{ij}\\}$ satisfy the normalization condition $\\sum_{i,j}\\h K_{ij}^\\dagger\\h K_{ij}=\\h I$.\n\nWhat is the natural quantification of the information in the setting of estimating an unknown state from the outcome of the quantum measurement?\nThe estimating process is characterized by a function $\\bm\\theta^{\\rm est}:I\\rightarrow \\Theta$ which is called an estimator.\nSince the outcome $i\\in I$ is a random variable, the estimator, which is calculated from the outcome, is also a random variable, and should be distributed around the true parameter $\\bm\\theta$.\nAccording to the classical Cram\\'er-Rao inequality, the variance-covariance matrix of the locally unbiased estimator has a lower bound which is determined only by a family of the probability distributions $\\{p_{\\bm\\theta}\\}_{\\bm\\theta\\in\\Theta}$:\n\\e{\n{\\rm Var}_{\\bm\\theta}[\\bm\\theta_{\\rm est}]\\ge (J^C_{\\bm\\theta})^{-1}.\n}\nHere, $J^C_{\\bm\\theta}$ is an $m\\times m$ matrix called the classical Fisher information whose elements are defined as\n\\e{\n[J^C_{\\bm\\theta}]_{ab}=\\sum_i p_{\\bm\\theta,i} \\pd{1}{\\log p_{\\bm\\theta,i}}{\\theta_a} \\pd{1}{\\log p_{\\bm\\theta,i}}{\\theta_b}. \\label{info}\n}\nTherefore, a measurement with a larger classical Fisher information allows us to estimate the state more accurately. \nWe define the classical Fisher information as the information obtained by the quantum measurement.\n\nThe disturbance can be evaluated as the loss of information on the parameter $\\bm\\theta$ that can be extracted from the quantum state.\nThe information on the parameter $\\bm\\theta$ that the quantum state potentially possesses can be quantified by the quantum Fisher information, which is defined as\n\\e{\n[J^Q_{\\bm\\theta}]_{ab}:=\\tr{\\pd{1}{\\h\\rho_{\\bm\\theta}}{\\theta_a}\\bm K^{-1}_{\\h\\rho_{\\bm\\theta}} \\pd{1}{\\h\\rho_{\\bm\\theta}}{\\theta_b}}.\n}\nHere, the superoperator $\\h K_{\\h\\rho}$ is defined as\n\\e{\n\\h K_{\\h\\rho}=\\bm R_{\\h\\rho}f(\\bm L_{\\h\\rho}\\bm R_{\\h\\rho}^{-1}),\n}\nwhere $f:(0,\\infty)\\rightarrow (0,\\infty)$ is an operator monotone function satisfying $f(1)=1$ and $\\bm R_{\\h\\rho} (\\bm L_{\\h\\rho})$ is the right (left)-multiplication of $\\h\\rho$:\n\\e{\n\\bm R_{\\h\\rho}(\\h A)=\\h A\\h\\rho, \\ \\bm L_{\\h\\rho}(\\h A)=\\h\\rho\\h A.\n}\nThe quantum Fisher information is the only metrics on the parameter space that monotonically decrease under an arbitrary completely positive and trace-preserving (CPTP) mapping $\\calE$~\\cite{Pe96}:\n\\e{\nJ^Q_{\\bm\\theta}(\\{\\h\\rho_{\\bm\\theta}\\})\\ge J^Q_{\\bm\\theta}(\\{\\calE(\\h\\rho_{\\bm\\theta})\\}).\n}\nFrom the monotonicity, the quantum Fisher information gives an upper bound of the classical Fisher information obtained by all the possible measurements, and the symmetric logarithmic derivative (SLD) Fisher information~\\cite{He68}, which corresponds to $f(x)=\\frac{1+x}{2}$, is known to be the least upper bound.\nTherefore, the quantum Fisher information, especially SLD Fisher information, can be interpreted as the information on the parameter $\\bm\\theta$ that can be extracted from the quantum state.\nWe define the disturbance of the measurement as \n\\e{\n\\Delta J^Q_{\\bm \\theta}:=J_{\\bm \\theta}^{Q}-\\sum_{i}p_ {\\bm\\theta,i} J_{i,\\bm \\theta}'^{Q},\\label{dist}\n}\nwhere $J_{\\bm \\theta}^{Q}$ and $J_{i,\\bm \\theta}'^{Q}$ are the quantum Fisher information of quantum statistical models $\\{\\h\\rho_{\\bm\\theta}\\}$ and $\\{\\h\\rho_{\\bm\\theta,i}\\}$, respectively.\nFor the sake of generality, we consider the disturbance using a general quantum Fisher information.\n\nThe information~\\eqref{info} and the disturbance~\\eqref{dist} satisfy the following inequality that shows a tradeoff relation between them:\n\\e{\nJ^C_{\\bm\\theta} \\le \\Delta J^Q_{\\bm\\theta}. \\label{i-d}\n}\nThis inequality means that if we perform a quantum measurement on an unknown state and extract information on the state, the state loses more intrinsic information.\nSince the inequality~\\eqref{i-d} is valid independently of the choice of the quantum Fisher information to define the disturbance, we obtain\n\\e{\nJ^C_{\\bm\\theta} \\le \\inf_Q \\Delta J^Q_{\\bm\\theta},\n}\nwhere the infimum is taken over all kinds of the quantum Fisher information to define the disturbance.\nWe note that it is nontrivial which quantum Fisher information gives the minimum disturbance, though the minimum and the maximum of the quantum Fisher information are known to be the SLD Fisher information and the real right logarithmic derivative Fisher information (real RLD), which corresponds to $f(x)=\\frac{2x}{x+1}$, respectively.\n\n\nThe proof of the inequality~\\eqref{i-d} is based on the monotonicity and the separating property of the quantum Fisher information.\nFor a given measurement $\\{\\h K_{ij}\\}$, we define a CPTP mapping $\\calE^{\\rm meas}$ as \n\\e{\n\\calE^{\\rm meas}(\\h\\rho)=\\bigoplus_i \\left( \\sum_j \\h K_{ij}\\h\\rho\\h K_{ij}^\\dagger \\right).\n}\nThen the quantum Fisher information of the quantum statistical model $\\{\\calE(\\h\\rho_{\\bm\\theta})\\}$ is equal to the sum of the classical Fisher information obtained by the measurement and the average quantum Fisher information of the post-measurement states as \n\\e{\nJ^Q_{\\bm\\theta}(\\{\\calE^{\\rm meas}(\\h\\rho_{\\bm\\theta})\\})=J^C_{\\bm\\theta}(\\{p_{\\bm\\theta}\\})+\\sum_{i}p_ {\\bm\\theta,i} J_{i,\\bm \\theta}'^{Q}, \\label{separating}\n}\nwhich we call the separating property (see Appendix~\\ref{proofA} for proof).\nBy applying the monotonicity under the CPTP mapping $\\calE^{\\rm meas}$, we obtain\n\\e{\nJ^Q_{\\bm\\theta}(\\{\\h\\rho_{\\bm\\theta}\\})&\\ge J^Q_{\\bm\\theta}(\\{\\calE^{\\rm meas}(\\h\\rho_{\\bm\\theta})\\}) \\nonumber \\\\\n&=J^C_{\\bm\\theta}(\\{p_{\\bm\\theta}\\})+\\sum_{i}p_ {\\bm\\theta,i} J_{i,\\bm \\theta}'^{Q},\n}\nwhich proves the inequality~\\eqref{i-d}.\n\n\\section{condition for equality}\nMeasurements that achieves the equality of the inequality~\\eqref{i-d} are efficient in a sense that they cause the minimum disturbance among those that extract a given amount of information.\nWhen we adopt the right logarithmic derivative (RLD) Fisher information~\\cite{YL73}, which corresponds to $f(x)=x$, to define the disturbance, a class of measurements called pure and reversible measurement achieves the equality of the inequality~\\eqref{i-d} (see Appendix~\\ref{proofB} for proof):\n\\e{\nJ^C_{\\bm\\theta}=\\Delta J^{\\rm RLD}_{\\bm\\theta} \\label{RLDequality}\n}\nHere, a measurement is called pure if the number of measurement operators is one for each measurement outcome so that the measurement operators can be written as $\\{\\h K_i\\}$, and is called reversible if each $\\h K_i$ has the inverse operator $\\h K_i^{-1}$~\\cite{UIN96}.\n\nAs an example, a measurement on a spin-1\/2 system proposed by Royer~\\cite{Ro94} is pure, whose measurement operators are \n\\e{\n\\h K_1&=\n\\begin{pmatrix}\n\\cos(\\theta\/2-\\sigma\/4)&0\\\\\n0&\\cos(\\theta\/2+\\sigma\/4)\n\\end{pmatrix},\\\\\n\\h K_2&=\n\\begin{pmatrix}\n\\sin(\\theta\/2-\\sigma\/4)&0\\\\\n0&\\sin(\\theta\/2+\\sigma\/4)\n\\end{pmatrix}.\n}\nThis measurement is reversible if $\\theta\/2\\pm\\sigma\/4\\ne n\\pi\/2$.\n\nIn fact, pure measurements are the least disturbing measurements in the following sense.\nSuppose that two measurements $\\{\\h K'_{ij}\\}$ and $\\{\\h K_i\\}$ give the same positive operator-valued measure (POVM)\n\\e{\n\\sum_j \\h K_{ij}'^\\dagger\\h K'_{ij}=\\h K_i^\\dagger \\h K_i, \\quad\\forall i\\in I,\n}\nand hence give the same probability distribution.\nThen, the pure measurement causes less disturbance, and gives the same amount of information:\n\\e{\n\\Delta J^Q_{\\bm\\theta}(\\{\\h K_i\\})&\\le \\Delta J^Q_{\\bm\\theta} (\\{\\h K'_{ij}\\}),\\\\\nJ^C_{\\bm\\theta} (\\{\\h K_i\\})&=J^C_{\\bm\\theta}(\\{\\h K'_{ij}\\}).\n}\n\n\n\\section{information-disturbance relation based on distinguishability}\nIn Ref.~\\cite{BL05}, by using the classical and quantum relative entropies\n\\e{\nS^C(p\\|q)&=\\sum_i p_i\\log\\left( \\frac{p_i}{q_i} \\right),\\\\\nS^Q(\\h\\rho\\|\\h\\sigma)&=\\tr{\\h\\rho(\\log\\h\\rho-\\log\\h\\sigma)},\\label{q.rel.}\n}\na similar inequality to~\\eqref{i-d} was derived:\n\\e{\nS^C(p\\|q)\\le S^Q(\\h\\rho\\|\\h\\sigma)-\\sum_i p_i S^Q(\\h\\rho_i\\|\\h\\sigma_i),\\label{rel.ent}\n}\nwhere $p,q$ and $\\h\\rho_i,\\h\\sigma_i$ are the probability distributions and the post-measurement states of a quantum measurement performed on the quantum states $\\h\\rho, \\h\\sigma$, respectively.\nSince the quantum relative entropy is a measure of the distinguishability of two quantum states~\\cite{HP91,ON02}, Eq.~\\eqref{rel.ent} can also be interpreted as a tradeoff relation between information and disturbance.\nIn particular, if we choose two similar states $\\h\\rho_{\\bm\\theta}$ and $\\h\\rho_{{\\bm\\theta}+{\\rm d}{\\bm\\theta}}$ as the arguments of the relative entropy, Eq.~\\eqref{rel.ent} reproduces the inequality~\\eqref{i-d} with the disturbance defined by the Bogoliubov-Kubo-Mori (BKM) Fisher information, which corresponds to $f(x)=\\frac{x-1}{\\log x}$.\n\nIn the following, we discuss the extension of the inequality~\\eqref{rel.ent} to general divergences.\nLet $D^C(\\cdot\\|\\cdot)$ be a divergence between two probability distributions, and $D^Q(\\cdot\\|\\cdot)$ be its quantum extension, i.e.,\nif two quantum states $\\h\\rho,\\h\\sigma$ commute and therefore are simultaneously diagonalizable as $\\h\\rho=\\sum_ip_i\\ket i\\bra i, \\h\\sigma=\\sum_iq_i\\ket i\\bra i$, we obtain\n\\e{\nD^Q(\\h\\rho\\|\\h\\sigma)=D^C(p\\|q).\n}\nWe note that quantum extensions of a divergence is not unique in general.\n\nLet us consider a condition for divergences $D^C(\\cdot\\|\\cdot),D^Q(\\cdot\\|\\cdot)$ to satisfy the information-disturbance tradeoff relation\n\\e{\nD^C(p\\|q)\\le D^Q(\\h\\rho\\|\\h\\sigma)-\\sum_i p_i D^Q(\\h\\rho_i\\|\\h\\sigma_i).\\label{divergence}\n}\nThe essential properties for the proof of the inequality~\\eqref{i-d} are the monotonicity and the separating property of the quantum Fisher information.\nIn the same way, we require these two properties for divergences:\n\\e{\nD^Q(\\h\\rho\\|\\h\\sigma)&\\ge D^Q(\\calE(\\h\\rho)\\|\\calE(\\h\\sigma)),\\\\\nD^Q(\\calE^{\\rm meas}(\\h\\rho)\\|\\calE^{\\rm meas}(\\h\\sigma))&=D^C(p\\|q)+\\sum_i p_iD^Q(\\h\\rho_i\\|\\h\\sigma_i).\n}\nIf we also require a continuity of $D^C(p\\| q)$ with respect to $p,q$, then it satisfies Hobson's five conditions that axiomatically characterize the classical relative entropy~\\cite{Ho69}.\nTherefore, the divergence with the monotonicity, the separating property, and the continuity must be consistent with the relative entropy at least for classical probability distributions:\n\\e{\nD^C(p\\|q)=S^C(p\\|q).\n}\nAs is shown in~\\cite{BL05}, the well-known quantum relative entropy satisfies the tradeoff relation~\\eqref{divergence} because it has the monotonicity and the separating property.\nAnother quantum extension of relative entropy proposed by Belavkin and Staszewski~\\cite{BS82},\n\\e{\nS^{\\rm BS}(\\h\\rho\\|\\h\\sigma)=\\tr{\\h\\rho\\log(\\h\\rho^{1\/2}\\h\\sigma^{-1}\\h\\rho^{1\/2})}, \\label{BS.rel.}\n}\nalso satisfies the inequality~\\eqref{divergence}.\nHere, $S^{\\rm BS}(\\cdot\\|\\cdot)$ is known to be maximal among all the possible quantum extensions of the classical relative entropy~\\cite{Ma13}.\nBy substituting $\\h\\rho_{\\bm\\theta}$ and $\\h\\rho_{{\\bm\\theta}+{\\rm d}{\\bm\\theta}}$ to the inequality~\\eqref{divergence}, we again obtain the inequality~\\eqref{i-d} with the disturbance defined by the real RLD Fisher information.\n\n\n\\section{conclusion}\nIn this paper, we have formulated the tradeoff relation between information and disturbance in quantum measurement in view of estimating parameters that characterize an unknown quantum state.\nThe information is defined as the classical Fisher information of the probability distributions of measurement outcomes, and the disturbance is defined as the average loss of the quantum Fisher information due to the backaction of the measurement.\nWe have shown the tradeoff relation~\\eqref{i-d} between them.\nWhen we use the RLD Fisher information, the equality of the inequality~\\eqref{i-d} is achieved by pure and reversible measurements.\nIn fact, pure measurements are the least disturbing among those that provide us with a given amount of information.\n\nWe have also discussed the necessary condition for divergences between two quantum states to satisfy a similar tradeoff relation~\\eqref{divergence}.\nIt is necessary for divergences to coincide with the relative entropy at least for classical probability distributions.\nIn addition to the well-known relative entropy, the maximum relative entropy also satisfies the tradeoff relation~\\eqref{divergence}, which reproduces the inequality~\\eqref{i-d} for the disturbance defined by the real RLD Fisher information.\nIf there are quantum extensions of the relative entropy that give an arbitrary quantum Fisher information, another systematic derivation of the inequality~\\eqref{i-d} should be possible.\n\\section*{acknowledgement}\nThis work was supported by\nKAKENHI Grant No. 26287088 from the Japan Society for the Promotion of Science, \na Grant-in-Aid for Scientific Research on Innovation Areas ``Topological Quantum Phenomena'' (KAKENHI Grant No. 22103005),\nthe Photon Frontier Network Program from MEXT of Japan,\nand the Mitsubishi Foundation.\nT. S. was supported by the Japan Society for the Promotion of Science through Program for Leading Graduate Schools (ALPS).\nY. K. acknowledges supports by Japan Society for the Promotion of Science (KAKENHI Grant No. 269905).\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{INTRODUCTION}\n\nHealth anxiety is a significant problem in our modern medical system \\cite{asmundson, taylor}. The belief that one's common and benign symptoms are explained by serious conditions may have several adverse effects such as quality-of-life reduction, incorrect medical treatment and inefficient allocation of medical resources. The web has been shown to be a significant factor in fostering such attitudes \\cite{whitehorvitzcyber, whitehorvitzhypo}. A recent study found that 35\\% of U.S. adults had used the Web to perform diagnosis of medical conditions either for themselves or on behalf of another person, and many ($>$50\\%) pursued professional medical attention concerning their online diagnosis \\cite{foxduggan}. Motivated by the popularity of online health search, we investigated how search engines might improve their health information offerings. We hypothesize that searchers are less likely to develop unrealistic beliefs when they are given unbiased and well-presented information about their medical state. Thus, we believe an ideal search engine should use queries, and available search histories, to extract medically relevant information about the individual as well as detect any health anxieties.\n\nLet us give a (negative) example of a possible medical search session. A user experiencing anxiety about a headache might first spend some time searching for information on a serious condition such as ``brain tumor'' and then switch to a symptom query about headache (a type of transition that prior work shows occurs frequently \\cite{cartright}). The user might choose the query wording ``severe headache explanations'' because of the subjective concern they experience at query time. The engine, registering the words ``severe'' and ``explanations'' as well as the phrase ``brain tumor\" present in the user's search history might compile a search engine result page (SERP) that is biased towards serious conditions. The user, viewing the SERP through the lens of their current health anxiety, may be attracted towards serious conditions in captions \\cite{whitehorvitzbias} and hence select a concerning page, heightening their anxiety further.\n\nIn this paper, we highlight a range of challenges and opportunities in working towards improving exploratory health search and thus hope to outline an agenda that frames this problem. Achieving this requires an understanding of both user and search engine behavior. Users play a role in the search process in two ways: their choice of query formulation as well their subjective consumption of information on the SERP and pages clicked on. This naturally led us to formulate three research questions:\n\\begin{itemize*}\n\\item [{\\bf Q1}] How does a user's subjective medical concern shape his or her choice of wording for medical queries?\n\\item [{\\bf Q2}] How does the search engine interpret the subjective medical concern and objective medical information expressed in the query as well as other measurable characteristics (such as medical search history) when compiling the SERP?\n\\item [{\\bf Q3}] How do users respond to the SERP, both in terms of the consumption of the information on the SERP as well as changes in future behavior caused by viewing the SERP and pages clicked on?\n\\end{itemize*}\n\nResults pertaining to these questions are found in Sections 4, 5, and 6 respectively. To answer them, we studied the search logs of 190 thousand consenting users of a major commercial Web search engine. Search logs are a valuable resource for studying information seeking in a naturalistic setting, and such data has been used by several studies to explore how searchers obtain medical information \\cite{whitehorvitzcyber, cartright}. The pipeline used to process this data and the features extracted for analysis are further described in Section 3. \n\nThe main contributions from our analysis are:\n\n\\begin{itemize*}\n\\item Revealing how certain users have specific preferences for certain query formulations (e.g. ``i have muscle pain'' vs. ``muscle pain'') which also has a significant effect on search results, and potentially health outcomes. \n\\item Finding evidence that users might not be swayed by concerning content appearing on SERPs as we might expect based on prior studies.\n\\item Quantifying the extent to which users with prior medical concern receive more concerning SERPs in response to health queries and choose the most concerning pages to click on, potentially leading to a vicious cycle of concern reinforcement. \n\\item Determining to how users directly and indirectly influence the level of concern expressed in the SERP they receive, both through query choice and other factors (e.g., personalization).\n\\end{itemize*}\n\nWe discuss these findings and their implications in Section 8 and conclude in Section 9.\n\n\\section{RELATED WORK}\n\nRelated research in this area falls into three main categories: health search behavior; the quality of online health content; and health anxiety and the impact of reviewing health content on the Web.\n\nThere continues to be interest in search and retrieval studies on expert and consumer populations in a variety of domains, typically conducted as laboratory studies of search behavior \\cite{bhavnani, hershhickam}. Benigeri and Pluye \\cite{benigeripluye} showed that exposing novices to complex medical terminology puts them at risk of harm from self-diagnosis and self-treatment. It is such consumer searching (rather than expert searching) that we focus on in the remainder of this section.\n\nSearch engine log data can complement laboratory studies, allowing search behavior to be analyzed at scale in a naturalistic setting and mined for a variety of purposes. Logs have been used to study how people search \\cite{jansen}, predict future search and browsing activity \\cite{lauhorvitz}, model future interests \\cite{dupretpiwowarski}, improve search engine quality \\cite{joachims}, and learn about the world \\cite{richardson}. Focusing on how people perform exploratory health searching, Cartright et al. \\cite{cartright} studied differences in search behaviors associated with diagnosis versus more general health-related information seeking. Ayers and Kronenfeld \\cite{ayerskronenfeld} explored changes in health behavior associated with Web usage, and found a positive correlation between it and the likelihood that a user will change their health behavior based on the content viewed. \n\nThe reliability of the information in search results is important in our study; unreliable information can drive anxiety. The quality of online healthcare information has been subject to recent scrutiny. Lewis \\cite{lewis} discussed the trend toward accessing information about health matters online and showed that young people are often skeptical consumers of Web-based health content. Eysenbach and Kohler \\cite{eysenbachkohler} studied users engaged in assigned Web search tasks. They found that the credibility of Web sites was important in the focus group setting, but that in practice, participants largely ignored the information source. Sillence and colleagues \\cite{sillence} studied the influence of design and content on the trust and mistrust of health sites. They found that aspects of design engendered mistrust, whereas the credibility of information and personalization of content engendered trust. \n\nThe medical community has studied the effects of health anxiety, including hypochondriasis \\cite{asmundson}, but not in Web search. Health anxiety is often maladaptive (i.e., out of proportion with the degree of medical risk) and amplified by a lack of attention to the source of their medical information \\cite{taylor, kring}. Such anxiety usually persists even after an evaluation by a physician and reassurance that concerns lack medical basis. A recent study showed that those whom self-identified as hypochondriacs searched more often for health information than average Web searchers \\cite{whitehorvitzhypo}. By estimating the level of health concern via long-term modeling of online behavior, search engines can better account for the effect that results may have and help mitigate health concerns. Our research makes progress in this area.\n\n\\begin{table*}\n\\centering\n\\begin{tabular}{>{\\centering\\arraybackslash}p{\\textwidth}} \nheadache, headaches, severe headache what do I do, which remedy for headache, headache top of head with back pain \\\\ \\hline\nheadache do I have a {\\it tumor}, {\\it headache rack}, my {\\it job} gives me a headache, headache {\\it national parks of california} \\\\ \\hline \na, low, be, he, sick, she, helps, standing, black, speech, male, between, acute, shaking, sensitive, bending, an, testing \\\\ \\hline\n\\end{tabular}\n\\captionsetup{justification=centering}\n\\caption{Examples of landmark queries (top), of non-landmark queries (middle; secondary topics and medical conditions are italic) and of admissible words we used to find potential landmark queries (bottom)}\n\\end{table*}\n\nSearchers may feel too overwhelmed by the information online to make an informed decision about their care\\cite{hart}. Performing self-diagnosis using search engines may expose users to potentially alarming content that can unduly raise their levels of health concern. White and Horvitz \\cite{whitehorvitzcyber} employed a log-based methodology to study escalations in medical concerns, a behavior they termed cyberchondria. Their work highlighted the potential influence of several biases of judgment demonstrated by people and search engines themselves, including base-rate neglect and availability. In a follow up study \\cite{whitehorvitzesc}, the same authors showed a link between the nature and structure of Web page content and the likelihood that users' concerns would escalate. They built a classifier to predict escalations associated with the review of content on Web pages (and we obtained that classifier for the research described in this paper). Others have also examined the effect of health search on user's affective state, showing that the frequency and placement of serious illnesses in captions for symptom searches increases the likelihood of negative emotional outcomes \\cite{lauckner}. Other research has shown that health-related Web usage has been linked with increased depression \\cite{bessiere}.\n\nMoving beyond the psychological impact of health search, researchers have also explored the connection between health concerns and healthcare utilization. In one study \\cite{whitehorvitzhui}, the authors estimated that users sought medical attention by identifying queries containing healthcare utilization intentions (HUIs) (e.g., [physicians in san jose 95113]). Eastin and Guinsler \\cite{eastinguinsler} showed that health anxiety moderated the relationship between health seeking and healthcare utilization. Baker and colleagues \\cite{baker} examined the prevalence of Web use for healthcare, and found that the influence of the Web on the utilization of healthcare is uncertain. The role of the Web in informing decisions about professional treatment needs to be better understood. One of our contributions in this paper is to demonstrate the potential effect of health-related result pages on future healthcare utilization.\n\nOur research extends previous work in a number of ways. By focusing on the first query pertaining to a particular symptom observed in a user history, we show that small differences in query formulation can reflect significant differences amongst health searchers and their health-related search outcomes. To date, no research has demonstrated the impact and insight afforded from analyzing such landmark queries and the behavior around them. To our knowledge, we are also the first to use the search logs to devise statistical experiments which allow us to quantify effects such as user response to medical search results amongst real user populations and provide evidence for causal relationships where possible. We believe that such an approach is necessary to formulate definitive implications for search engine design as well as measuring search engine performance. \n\n\\section{STUDY}\nWe describe various aspects of the study that we performed, including the data, the features extracted, and the statistical methods used for analysis.\n\n\\subsection{Log Data}\n\nTo perform our study we used the anonymized search logs of a popular search engine. Users of this search engine give consent to having information about their queries stored via the terms of use of the engine. During this study, we focused on medical queries related to headache, as it is among the most common health concerns \\cite{nielsen}. We use the phrase {\\it headache query} to refer to queries that contain the substring ``headache'' and occurred during the six month period from September 2012 to February 2013. Amongst those, we call a {\\it landmark query} a query that shows an intent to explore the symptom ``headache'', that does not already contain a possible explanation for headache (e.g., migraine, tumor) and that is not otherwise off-topic. We found these landmark queries by manually assessing frequent headache queries and creating a list of 682 ``admissible'' terms that we believed could occur in landmark queries. We then compiled all queries that exclusively contain terms from that list into a dataset. Examples of landmark queries, non-landmark queries, and admissible words can be found in Table 1. From manual inspection, we concluded with confidence that the dataset captured over 50\\% of all landmark queries present in the logs and that over 95\\% of captured queries are proper landmark queries, the rest being headache queries which contain a significant secondary topic. We then excluded all but the first landmark query instance for each user, ensuring that each user only appears once in the dataset. Overall, our dataset contains over 50,000 unique queries and over 190,000 query instances \/ users.\n\nWe focus on headache since it a common medical symptom (e.g., over 95\\% of adults report experiencing headaches in their lifetime \\cite{rasmussen}) and there are a variety of serious and benign explanations (from caffeine withdrawal to cerebral aneurism), facilitating a rich analysis of content, behavior, and concern. While we believe that headache searching is sufficiently rich and frequent to warrant its own study, investigating queries related to symptoms other than headache could solidify our findings. A large-scale analysis similar to that reported here, but focused on multiple symptoms, is an interesting and important area for future work.\n\n\\subsection{Features}\n\nFor each landmark query in our dataset, we generated features. To frame our three research questions, we modeled the search process around a landmark query as five separate stages: (i) The user's search behavior prior to the landmark query, (ii) the user's choice of wording of the landmark query, (iii) the SERP returned to the user by the search engine, (iv) the user's decision which pages to click on (if any), and (v) the user's search behavior after the landmark query. Our research questions ask about the relationship of these 5 stages. Hence, the features we extracted come in five groups.\n\n\\begin{table}\n\\centering\n\\begin{tabular}{lm{2.3in}cc} \nName&Description \\\\ \\hline\npasttopserious&Number of medical queries containing the most frequent serious condition amongst all serious conditions present \\\\ \\hline\npasttopbenign&Number of medical queries containing the most frequent benign condition amongst all benign conditions present \\\\ \\hline\npastdiffserious&Number of distinct serious conditions in medical queries \\\\ \\hline\npastdiffbenign&Number of distinct benign conditions in medical queries\\\\ \\hline\npastmedical&Number of medical queries \\\\ \\hline\npastheadache&Number of medical queries containing ``headache'' or a condition that is an explanation for headache \\\\ \\hline\npasthui&Number of HUI queries \\\\ \\hline\npasthuiclicked&Number of HUI queries where at least one page on the SERP was clicked \\\\ \\hline\n\\end{tabular}\n\\caption{User Search Behavior Features}\n\\end{table}\n\n\\begin{table*}\n\\centering\n\\begin{tabular}{lllc} \nName & Description (query ..) & Example & Frequency\\\\ \\hline\naudience & .. is about specific population &``headache in adults'' & 3.6\\%\\\\ \nfiller & .. contains a filler such as 'a' or 'and' &``headache and cough'' & 3.6\\%\\\\ \ngoal & .. contains a specific search goal & ``definition of headache'' & 26.3\\%\\\\ \ngoal:condition & .. indicates the goal of diagnosis & ``reasons for headache'' & 8.2\\%\\\\ \ngoal:symptom & .. indicates the goal of related symptoms & ``headache symptoms'' & 2.1\\%\\\\ \ngoal:treatment & .. indicates the goal of treatment & ``headache cure'' & 14.7\\%\\\\ \ngoal:medication & .. indicates the goal of treatment through medication & ``headache pills'' & 2.4\\%\\\\ \ngoal:alternative & .. indicates the goal of alternative treatment & ``natural headache remedy'' & 5.4\\%\\\\ \nsymptomhypothesis & .. states another symptom as cause & ``headache caused by back pain'' & 1.7\\%\\\\ \neventhypothesis & .. relates the headache to a life event & ``headache after hitting head'' & 3.3\\%\\\\ \nduration & .. specifies a duration for the headache & ``chronic headache'' & 10.9\\%\\\\ \nintensity & .. specifies that the headache is strong & ``severe headache'' & 5.0\\%\\\\ \nlocation & .. specifies the location of the headache & ``headache left side of head'' & 17.3\\%\\\\ \npronoun & .. contains a pronoun & ``i have headache'' & 5.8\\%\\\\ \nkindofheadache & .. specifies the kind of pain & ``stabbing headache'' & 2.5\\%\\\\ \nothersymptom & .. specifies an additional symptom & ``headache reflux'' & 27.9\\%\\\\ \ntriggered & .. indicates a headache trigger & ``headache when bending over'' & 1.4\\%\\\\ \ntimeofday & .. indicates a daily pattern & ``headache in afternoon'' & 3.2\\%\\\\ \nopenquestion & .. is phrased as an open question & ``what to do about headache'' & 11.1\\%\\\\ \n\\end{tabular}\n\\caption{High-level {\\it QueryFormulation} features. Note that queries may have multiple features, such as ``severe headache on top of head and cough'', so the percentages in the far-right column do not sum to 100\\%.}\n\\end{table*}\n\n{\\it BeforeSearching}. This group of features describes the level of medical searching before the landmark query (stage (i)). For each query in the user's search log, we first extracted whether that query was of a medical nature. For this, we used a proprietary classifier. From manual inspection, we concluded with confidence that its Type I and II errors are $< 0.1$. Queries so classified as medical will be called {\\it medical queries}. Secondly, we extracted phrases present in the query that describe medical conditions, such as ``common cold'' or ``cerebral aneurism''. The list of phrases we considered was based on the International Classification of Diseases 10th Edition (ICD-10) published by the World Health Organization as well as the U.S. National Library of Medicine's PubMed service and other Web-based medical resources. We also used manually curated synonyms from online dictionaries and standard grammatical inflections to increase coverage. For more information see the approach used by \\cite{whitehorvitzcyber}. The list was also separated into benign and serious medical conditions and we determined which conditions are possible explanations for headache. Thirdly, we extracted phrases that indicate the query is linked to an intention of real-world healthcare utilization (HUI), such as ``emergency care'' or ``hepatologist''. We call those queries {\\it HUI queries}. Finally, the individual-query-level features were aggregated over a time window just before the landmark query to form 8 distinct features, which are shown in Table 2. In our experiments, we considered five different aggregation windows: 1 hour, 1 day, 1 week, 30 days and 90 days. All experiments involving {\\it BeforeSearching} features were replicated for each aggregation window. \n\nWe believe that intense medical searching may be a sign of health concerns or anxieties. White and Horvitz \\cite{whitehorvitzhypo} extensively demonstrated that users believing their symptoms may be explained by a serious condition conduct longer medical search sessions and do so more frequently. Of course, there are many potential reasons for increased medical search such as different web search habits, random noise or even a recent visit to a physician \\cite{whitehorvitzhui}. Overall, we believe that it makes sense to view {\\it BeforeSearching} in light of possible medical concern experienced. However, even if there are significant other factors at play, we believe that the phenomenon of health search intensity remains interesting.\n\n{\\it QueryFormulation}. The choice of wording for the query (stage (ii)) was modeled, firstly, using 19 high-level features. We arrived at these by manually inspecting admissible words and landmark queries (such as in Table 1) and noting the most important high-level ideas expressed through them. We believe that those 19 features capture a significant portion of the semantic variation within queries. The features are shown in Table 3. They offer a useful characterization of the broad range of different types of search intent associated with headache-related queries. Four of these features ({\\it othersymptom}, {\\it duration}, {\\it intensity} and {\\it location}) were further divided by which key phrase was used to express this feature, yielding an additional 117 low-level features (examples are shown in the graph annotation of Figure 1.2-1.5). All feature extraction functions are based on substring matches joined by logical operators. From manual inspection, we concluded with confidence that the Type I and II errors of the extractors of all these features was low. All features in this category are binary.\n\n\\begin{figure*}\n\\centering\n\\includegraphics[scale=1.8]{F1p1.pdf}\n\\includegraphics[scale=1.8]{F1p2.pdf}\n\\includegraphics[scale=1.8]{F1p3.pdf}\n\\caption{Association between query formulation and search behavior features. From left to right: {\\bf {\\color{red} Red}}: {\\it pastmedical} (aggregation window: 90 days); {\\bf {\\color{blue} Blue}}: {\\it pastmedical} (window: 1 hour); {\\bf {\\color{Brown} Brown}}: {\\it futuremedical} (window: 1 hour); {\\bf Black}: {\\it futuremedical} (window: 90 days). All bars show the relative change of medical search intensity of users whose landmark query has a certain formulation relative to the global mean. Error bars are 95\\% confidence intervals.}\n\\end{figure*}\n\n{\\it SERPConcern}. We scored the level of medical concern expressed in each of the Top 8 pages on the SERP (stage (iii)) using a logistic regression classifier designed to predict searching for serious conditions and shown in \\cite{whitehorvitzesc} to have significant predictive power. The classifier was graciously provided to us by the authors for use in our study. Note that during the time period analyzed in this study, the search engine only returned eight results on a large number of SERPs so we disregarded possible further results. It has been shown that users rarely click below position 8, including for health queries \\cite{whitehorvitzbias}. The classifier is based on page features from the URL and HTML content similar to those shown in Table 2. It also includes features that attempt to measure the ``respectability'' of the page (e.g., expressed through the {\\it Health on the Net Foundation} certificate \\cite{hon}) to capture the effect of this on the user. It is then trained to discriminate between pages that lead to serious condition searches (concerning) and pages that lead to benign condition searches (non-concerning). The concern score of an arbitrary page is then the inner product between its feature vector and the learnt weight vector. Finally, we take the weighted sum of the 8 scores for the individual pages on any SERP to produce a single feature value for the full SERP. Note that even though we consider pages leading to benign condition searches less concerning than those leading to serious condition searches, throughout this study, we still consider benign condition searching as an indicator for medical concern experienced by the user, albeit less strong than serious condition searching.\n\n{\\it ClickFeatures}. To measure the users click behavior (stage (iv)), we record whether the user clicked a page on the SERP. If the user did, we also recorded the concern score of the page clicked (as described in the last paragraph) as well as the position of the clicked page on the SERP. We call these three features {\\it hasclicked}, {\\it clickconcern} and {\\it clickposition}. (Users that did not click are excluded from analysis involving features {\\it clickconcern} and {\\it clickposition}.)\n\n{\\it AfterSearching}. The same 8 features as {\\it BeforeSearching}, but aggregated over a window just after the landmark query (stage (v)). In the name, we replace ``past'' with ``future'' (e.g. one feature is called {\\it futuremedical}).\n\n\\subsection{Statistical Methodology}\n\n\nLet $X$ be a dataset where each entry corresponds to a user \/ landmark query (either our full dataset of 190.000 entries or a subset of it.). Write $x_1, .., x_N$ for the data points and $x^1_n, .., x^d_n$ for the components \/ feature values of data point $x_n$. Throughout our analysis, we wish to measure whether there is a significant association between two feature values $i$ and $j$, say between {\\it pastmedical} and {\\it SERPConcern}. By this we mean that the p-value of a suitable independence test on the two variables is low. To measure this simply and robustly, we will choose a threshold $t^i$ and split the data set into two subpopulations $X_<$ and $X_>$ such that for all $x_n$, $n \\in \\{1, .., N\\}$, we have $x_n \\in X_< \\iff x_n^i < t^i$. We also call the two subpopulations the {\\it lower bucket} and {\\it upper bucket} respectively. Then, we either perform a two-sample t-test on the population means of $X_<$ and $X_>$ (Figure 3.3) or we consider a 95\\%-confidence interval around the mean of the smaller bucket if it is significantly smaller than the other bucket (Figures 1 and 2). (In practice, all our subpopulations are large enough to warrant the use of gaussian confidence intervals \/ tests.) A statistical association between features $i$ and $j$ is a symmetric relation, we may choose to split on either feature and compare the means of the other, based on convenience of presentation.\n\nSeveral times, we will encounter a more challenging case where we want to answer the question whether two feature values $i$ and $j$ are associated while controlling for a third feature value $k$. We do this by first dividing the dataset into many subpopulations according to the exact value of feature $k$ to obtain $X^{v^k_1}$, .., $X^{v^k_{N^k}}$, where $v^k_1, .., v^k_{N^k}$ are the values $x_n^k$ can take. Then, we split each of these subpopulations further according to thresholds on feature $i$ as before. Hence, we obtain $X_<^{v^k_1}$, .., $X_<^{v^k_{N^k}}$ and $X_>^{v^k_1}$, .., $X_>^{v^k_{N^k}}$. Because it is nontrivial to jointly compare this potentially large number of buckets, we adopt the following 2-step procedure. \n\nFirst, we take the union of all lower and upper buckets respectively and perform a two-sample t-test as before. For this to control for feature $k$, each lower bucket must be of the same size as its corresponding upper bucket. If, for some $X^{v^k_n}$ there is no threshold that achieves this, we randomly remove data points whose feature value is equal to the median across $X^{v^k_n}$ until this is possible. While this significantly reduces the number of data points entering the analysis, we do not believe this threatens validity or generalizability.\n\nIn the second step, we first individually compare each lower bucket to its corresponding upper bucket. Because many of these buckets are small (e.g. of size 1), we use the Mann-Whitney U statistic to obtain a p-value for each pair of buckets. Then, we aggregate all of those p-values by Stouffer's Z-score method where each p-value is weighted according to the size of its respective buckets. \n\nFinally, we wish to combine p-values obtained in both steps to either accept or reject the null-hypothesis at a given significance level. This is achieved by taking the minimum of both values and multiplying by 2. The cumulative distribution function of that combined quantity is below that of the uniform distribution under the null hypothesis and is thus at least as conservative in rejecting the null as each individual p-value, while gaining a lot of the statistical power of both tests. We quote this value in Figure 3.1-3.2 and 3.4-3.7.\n\n\\section{QUERY FORMULATION}\n\n\\begin{figure*}\n\\centering\n\\includegraphics[scale=1.8]{F2v2.pdf}\n\\caption{Association between query formulation features and {\\it SERPConcern}. All bars show the difference in {\\it SERPConcern} of users whose landmark query has a certain formulation relative to the global mean, measured in standard deviations of the global distribution.}\n\\end{figure*}\n\nWe began our analysis by investigating how users with different levels of {\\it BeforeSearching} phrase their queries. Figure 1.1 shows the mean value of {\\it pastmedical} for each high-level {\\it QueryFormulation} feature relative to the population mean, both aggregated over 1 hour and 90 days. Error bars indicate 95\\% confidence intervals. We chose the feature {\\it pastmedical} to represent {\\it BeforeSearching} because it is the least sparse and hence has the smallest confidence interval.\n\nTo our surprise, we found that most {\\it QueryFormulation} features are highly discriminative, i.e. users choosing queries with certain features have conducted significantly more medical searching than users choosing queries that have different features. In fact, every one of these 19 features is associated with a significant change in prior medical search activity during the hour before the landmark query. This effect could be explained by users near the beginning of their medical search process choosing different query patterns compared to users near the end of their medical search process. However, most {\\it QueryFormulation} features are still discriminative when aggregated over a 90 day window (red bars), and we found that the majority of medical searches in that window do not occur immediately prior to the landmark query. Figure 1.1 also shows the feature {\\it futuremedical}. We find that most {\\it QueryFormulation} features have the same association with past and future search activity, even over a 90 day window. This is evidence that these {\\it QueryFormulation} features characterize users and that heavy medical searchers prefer certain formulations compared to other users in a consistent fashion over the long term.\n \nOver the 1 hour window, users entering an additional symptom in their landmark query show the highest level of search activity. This might be because experiencing a larger number of symptoms causes the user to want to find information about each individual symptom, thus increasing the search need. Surprisingly, {\\it intensity} is not one of the strongest predictors of increased searching even though it appears to be the most intuitive indicator of increased concern. Also, the features {\\it openquestion} and {\\it pronoun} indicate less prior searching. Hence, the presence of sentence-like structures in queries is potentially associated with lower user health concern. It is difficult to intuitively interpret most of the {\\it QueryFormulation} features. Figure 1.2-1.5 shows low-level {\\it QueryFormulation} features relating to high-level features {\\it duration}, {\\it intensity} and {\\it othersymptom}. Unfortunately, the sparseness of {\\it pastmedical} leads to large margins of error, the exact size of which are also difficult to determine. Nonetheless, certain features which are more common and thus have lower margins or error such as ``constant'' or ``terrible'' are as discriminative as high-level features, suggesting they are useful for making inferences about the user. For example, our analysis suggests the very counterintuitive conclusion that users searching for ``headache everyday'' experience a different level of concern than users searching for ``daily headache'' (see difference in feature {\\it pastmedical} between {\\it daily} and {\\it everyday} (blue bar) in Figure 1.2). Similarly, different symptoms are associated with their own level of past search activity. We chose to display 10 relatively concerning symptoms (Figure 1.4) and 10 relatively benign symptoms (Figure 1.5). The choice was made using the findings of a separate crowdsourcing study that we omit from the paper due to space reasons. In that study, many participants were asked to estimate the level of medical concern associated with a set of symptoms. Even though we do not present this study, the difference between the two groups of symptoms is intuitively clear. Surprisingly, we do not see a clear trend that users under Figure 1.4 have searched more, which weakens the hypothesis that objective medical state is linked to search activity.\n\nIn summary, we find that both high-level {\\it QueryFormulation} features as well as individual word choices reveal information about the searcher, which is not necessarily expected. While some {\\it QueryFormulation} features are interpretable, more work is necessary to understand their precise meaning. Also, more data is needed to better study their effect on rarer search events such as HUI queries. \n\n\\section{EFFECT ON SERP}\n\n\n\nWhen personalization is employed by search engines a ``filter bubble\" can be created whereby only supporting information is retrieved \\cite{pariser}. As highlighted in the earlier example, this can be problematic in the case of health searching. In this section, we investigate in what ways the user influences the content of the SERP and the medical concern expressed therein. This question can be separated into two aspects: How do SERPs vary for a given query and how do SERPs vary between queries. \n\n\\subsection{Same Query, Different SERPs}\nTo determine the impact of the user on SERPs within each query, we first investigated how diverse those SERPs were to begin with. We analyzed the composition of SERPs of the 29 most frequent queries from our data set. (Each of these occurred at least 500 times.) We found that on average, 92\\% of SERPs contain the same top result. For example, the page ``www.thefreedictionary.com\/headache'' appears as the top result for the query ``headache definition'' for almost every user. The three most common top results together covered over 99\\% of SERPs. If we consider the Top 8 results on the SERP, we find that eight specific pages are enough to account for 61\\% of all results. Hence, most SERPs returned for a given query are highly similar. We do see considerably more diversity when we consider the ordering of pages on the SERP. Hence, we conclude that factors such as time of day, user location and user personalization may shuffle the ordering of pages, especially in the lower half of the SERP, but do not have the power to promote completely different pages to the top the majority of the time (one notable exception to this is personal navigation \\cite{teevan}, which agressively promotes pages that an individual visits multiple times, but the coverage of this approach is small). Since user outcomes are driven chiefly by top results, we thus expect the impact of non-query factors on user outcomes to be limited.\n\n\n\nNonetheless, we investigated the impact of prior search activity on the SERP by measuring the association between {\\it BeforeSearching} features and {\\it SERPConcern} while controlling for query choice. We measure significance as described in section 3.3. We split on {\\it SERPConcern} and compared the empirical means of the {\\it BeforeSearching} features. Figure 3.1 shows the percentage difference between the mean of the union of the upper buckets and the mean of the union of the lower buckets, aggregated over a 24 hour window before the landmark query (the largest window where significance was obtained). The p-value is shown above the bars. We only show {\\it BeforeSearching} features that were significant. \n\nWe see that users who received concerning SERPs relative to other users entering the same query searched for 3\\% more serious conditions and 2\\% more for the most frequent serious conditions in the 24 hours before the landmark query than than those receiving less concerning SERPs. So, there is an impact of prior searching on the SERP, albeit, as expected, a small one. The difference between the buckets becomes much larger when we consider only users who receive the 10\\% most and least concerning SERPs within each query. Figure 3.2 shows that users receiving especially concerning SERPs search for 12\\% more serious conditions, search 18\\% more for their most frequent serious condition and conduct 10\\% more medical searches overall in the 24 hours before the landmark query. If it was the case that prior searching indicates heightened concern, then search engines present significantly more concerning search results to already concerned users, which may be undesirable.\n\n\\subsection{Different Queries, Different SERPs}\n\n\\begin{figure*}\n\\centering\n\\includegraphics[scale=1.8]{F3p1.pdf}\n\\includegraphics[scale=1.8]{F3p2.pdf}\n\\caption{Association between user search features and features related to the landmark query. Analysis was conducted as described in section 3.3. In each case, we show the ratio between the search behavior feature value across all upper buckets and across all lower buckets. Hence, positive values imply a positive association between features and vice versa. The p-value is noted next to the bar. We do not show results that were not significant, except in Figure 3.7, were only {\\it futurehuiclicked} was significant. The effective sample size refers to the total number of individuals \/ datapoints included in the analysis. This varies for three reasons. (1) Some features are unavailable for some users. For example, in Figure 3.6, we cannot include users that did not click on the SERP. (2) In some Figures, we only include a specific subset of users. For example, in Figure 3.2, we only include users receiving the 10\\% most \/ least concerning SERPs for each query. (3) We have to subsample for statistical reasons (see section 3.3).}\n\\end{figure*}\n\nNow, we turn our attention to how SERPs vary across queries. As a starting point, we tested the association between {\\it BeforeSearching} and {\\it SERPConcern} (without controlling for query choice). Results are shown in Figure 3.3. {\\it BeforeSearching} values shown were aggregated over a one hour window. This window size yielded the most significant results, but significance was preserved over all window sizes up to 90 days. Interestingly, users receiving concerning SERPs are now significantly less likely to engage in medical searching before the landmark query. They search for 14\\% less benign conditions and 17\\% less for their most frequent benign condition. Since these values (and p-values) are much larger than in Figure 3.1, we conclude that the choice of query is the cause for the majority of this effect. At face value, this effect may be either desirable or counterproductive depending on whether {\\it BeforeSearching} is an indicator of subjective concern or objective health state.\n\nIn Figure 2.1, we break down {\\it SERPConcern} by high-level {\\it QueryFormulation} features. Error bars are much smaller than those for {\\it pastmedical} because the SERP scoring function is not sparse. We see that including an additional symptom in the landmark query (feature: {\\it othersymptom}) significantly lowers the level of concern in the SERP. This might be because most of these symptoms are not indicative of a serious brain condition (e.g. cough, stomach pain, hot flashes etc.), thus providing evidence to the search engine that such a condition might not be the underlying reason of the user's health state. We note that this association is the opposite of the association between {\\it othersymptom} and {\\it pastmedical}. We hypothesize that this effect might be responsible for the negative association in Figure 3.3. Indeed, if we exclude searches that include additional symptoms, all negative associations in Figure 3.3 disappear and we obtain significant positive associations between {\\it SERPConcern} and prior serious condition searching. \n\nIn contrast, searching specifically for conditions that explain headache (feature: {\\it goal:condition}) yields the most concerning pages. This is logical given what makes a page most concerning is content about (serious) conditions. Again, it is difficult to interpret the meaning of most of the {\\it QueryFormulation} features intuitively, but each feature is highly discriminative with respect to {\\it SerpConcern} and is thus warranting further study.\n\nFigure 2.2-2.5 breaks down SERPConcern by low-level {\\it QueryFormulation} feature. It appears that search engines do a decent job at ranking different symptoms based on how medically serious they are. {\\it SERPConcern} is overall much higher amongst relatively concerning symptoms versus relatively benign symptoms, showing that the search engine does respond to this factor. Furthermore, we see a form of consistency in the search engine in that queries with additional symptoms consistently receive below average concern scores. We compared this against different key phrases describing the feature {\\it location} (see Table 3) such as ``above eye'', ``temporal'' and ``base of skull''. If we model the {\\it SERPConcern} value associated with each additional symptom as well as each location descriptor as normally distributed, then the difference in mean between the two normals is highly significant ($p < 10^{-6}$).\n\nNonetheless, we still see that there is a considerable amount of unexplained variation in {\\it SERPConcern} values. For example, a user entering ``bad headache'' will receive a less concerning SERP than a user entering ``terrible headache''. We hypothesize that it is unlikely that this difference is always due to hidden semantics, but may often be attributable to random noise. This suggests that there is still significant scope for search engines to better reflect the medical content of queries and become more consistent. One caveat is that some of this noise might be caused by our classifier that is used to assign {\\it SERPConcern} values. The presence of this noise implies that different user populations which have a preference for certain query formulations may be inadvertently led down completely different page trails to different health outcomes.\n\n\\subsection{Summary}\n\nIn summary, we find that search engines do seem to show the ability to return less concerning SERPs for benign symptoms in queries as opposed to more serious symptoms. However, there is significant work to be done to achieve a state where search results reflect the medical information given in the query while being robust to irrelevant nuances. Also, we find that past user searching does have a direct impact in making SERPs more concerning, which might not be desirable. Further analysis might discover the exact cause.\n\n\\section{RESPONSE OF USER TO SERP}\n\nIn this section, we investigate how users respond to the SERP returned in response to the landmark query. We phrase this as two sub-problems: How do users with different levels of prior searching interpret the SERP by making click choices and what impact does the SERP have in altering the behavior of the user. \n\n\\subsection{Impact on Click Behavior}\n\nFirst, we looked at the impact of {\\it BeforeSearching} on click decisions. For this, we measured the association between {\\it BeforeSearching} and {\\it ClickFeatures}. However, to capture the impact of user predisposition, we can only compare users against each other who not only entered the same query, but also received identical SERPs, to control for those two confounders. Unfortunately, this means we cannot include users in our study that received a SERP none else received, which significantly reduces the size of our effective dataset to lie between 20,000 and 50,000 and users. We split each subpopulation based on high \/ low values of {\\it hasclicked}, {\\it clickconcern} and {\\it clickposition} and compared the means of {\\it BeforeSearching} features. Results are shown in Figures 3.4 (window: 30 days), 3.5 (window: 30 days) and 3.6 (window: 24 hours) respectively. \n\nWe found that users who select more concerning pages on any given SERP are significantly more likely to have conducted medical searches in the 30 days before the landmark query. They have searched for 19\\% more serious conditions over this time period. This is quite large given the size of the aggregation window and obfuscating factors such as the limited amount of information available about a page in a SERP caption, the ad hoc nature of a click decision in general and the overall preference for pages ranked near the top. This suggests that people are selectively seeking information and concerned users reinforce their opinion by focusing on concerning content on the SERP. Selective exposure to information has been studied in detail in the psychology community \\cite{frey}. The fact that the significance of these results is highest for a large aggregation window illustrates a user's page preference is formed over the medium term, suggesting it is an attitude rather than a momentary state. Furthermore, this result points to past medical search activity as a good proxy for level of health concern, making our previous analyses more meaningful.\n\nWe found that users with more prior health searching are more likely to click lower positions on the SERP and are less likely to click overall. This might be because users who have searched about similar topics before are likely to look for specific kinds of information and are thus more likely to reject pages as unsuitable, leading them further down the SERP and ultimately to abandon more of their queries. Connections between topic familiarity and search behavior have been noted in previous work \\cite{kelly}. Interestingly, {\\it huiclicked} breaks that trend and is positively associated with clicking on the SERP. This might be because {\\it huiclicked} is by definition associated with a user's general disposition to click on a SERP, which in turn affects the probability of the user clicking in response to the landmark query. \n\n\\subsection{Impact on Future Behavior}\n\nNow we turn to measuring the impact that the SERP has on the state of the user as measured by {\\it AfterSearching}. Previous research (e.g., survey responses in \\cite{whitehorvitzcyber,foxduggan}) showed that online content can have a direct impact on people's healthcare utilization decisions. It also has been shown that the content of pages viewed can be used to predict future serious condition searches \\cite{whitehorvitzesc}. We believe that finding evidence in logs that different SERPs cause users to respond differently would strengthen the motivation for improving search engine performance during health searches. \n\nBefore proceeding, we must point out that this task is quite difficult. We have already shown that users have a myriad of significant predisposition which shape the SERP. Hence, it is impossible to tell whether any change in user behavior after the landmark query was really caused by the SERP or is the result of a predisposition.\n\nOne way to mitigate this is, again, control for query choice. Unfortunately, the group of users receiving concerning SERPs within each query is very different from the group of users receiving concerning SERPs overall, as within-query variation of {\\it SERPConcern} is much smaller than between-query variation (see Section 5). Hence, the true impact of the SERP on users is likely much larger than is measurable in this experiment.\n\nSplitting on {\\it SerpConcern} and comparing the means of {\\it AfterSearching} features yielded no significant differences over any aggregation window. However, looking at the top 10\\% vs. bottom 10\\% of SERPs within each query yielded surprising results. They are shown in Figure 3.7 (window: 1 hour). Users receiving especially concerning SERPs perform 20\\% less HUI queries in the hour following the landmark query when compared to their peers who entered the same query but received an especially unconcerning SERP. Additionally, these users had a larger empirical probability of expressing HUI in the preceding hour (non-significant), amplifying the drop. The ratio of HUI queries between users receiving concerning vs. non-concerning SERPs changes from +20\\% to -20\\% from past to future. To solidify this result, we reran our experiment under exclusion of users who had performed any medical searching over the last 1 hour, 24 hours, etc. We found that the negative association of {\\it futurehui} and {\\it futurehuiclicked} with {\\it SERPConcern} remained significant even when excluding users who had performed no medical searching in the preceding 30 days. (Note that the more users we exclude, the less data we have and the more difficult it is to achieve significance.) This is a surprising result that warrants further investigation.\n\nEven though there is no causal argument to be made, we were still interested in the outright association of {\\it SERPConcern} and {\\it AfterSearching}. We made the interesting observation that the association of {\\it topserious} and {\\it diffserous} actually decreased from past to future. For example, users receiving concerning SERPs issued 2\\% more searches for their most frequent serious condition during the hour before the landmark query, but 1\\% less afterwards (both non-significant). On the contrary, users receiving concerning SERPs searched 17\\% less for their most frequent benign conditions in the hour before the landmark query and only 8\\% less in the hour after. These results are difficult to interpret. However, they do provide some evidence that the impact of the SERP on the user's concern might be more subtle than expected.\n\n\\subsection{Summary}\nIn summary, we found that concerned users are significantly more likely to click concerning pages which may be a serious problem. Additionally, we found only weak evidence that concerning SERPs cause an increase in user's health concerns. To the contrary, users receiving concerning SERPs may be less likely to seek HUI in the short term. Both of these effects need to be investigated further.\n\n\\section{DISCUSSION AND IMPLICATIONS}\n\nIn this section, we summarize our findings and discuss implications, limitations and opportunities for future research. Our main findings are as follows:\n\n\\begin{itemize*}\n\\item Innocuous details in query formulation can hold characterizing information about the search engine user.\n\\item While the search engine does respond to general trends in query formulation sensibly, there is a lot of variation in SERPs for different but similar word choices. \n\\item While users with a history of medical searching are significantly more likely to pick out especially concerning content, the search engine also serves those users more concerning SERPs to begin with.\n\\item There was only weak evidence for the effect of intensification of health concerns through concerning SERPs, but some evidence that concerning SERPs might reduce HUI queries and hence real-world health seeking in the short term.\n\\end{itemize*}\n\nConsidering this, we believe that the most important implication for improving medical search results from this study is clearly highlighting the need for a further rigorous quantification of the extent to which concern in SERPs influences users' levels of health concern. Utlimately, through methodologies to study consenting user cohorts in detail, we can move beyond our focus on health concerns to target health anxieties directly.\n\nWe also believe there is significant potential in refining our understanding of the meaning of both landmark query formulation and general search behavior for extracting hidden information about the user. In this paper, we have shown how these might be used to better infer levels of health concern. In this way, we might be able to improve search engine performance, for example by promoting more medically trusted pages that discuss a wide variety of possible causes objectively if we suspect health anxiety in the searcher. However, a lot more can be done. For example, word choices might reveal the age of the user \\cite{torres} or their level of domain expertise \\cite{hembrooke}, which in turn would have implications for the medical meaning of symptoms. We also have not yet considered formulation nuances of queries other than the landmark query, which might hold much richer information. The need for such analysis to adjust for anxiety is exemplified by our finding that users who conduct more medical searching have a preference for concerning content, which might indicate a cycle of self-reinforcing concern.\n\nOther questions for future research include: In how far is frequent medical searching indicative of the user's actual state of health vs. subjective concern? How can we effectively integrate the results of a possible algorithm that learns searchers' true medical situation into a standard retrieval model without ``overmanaging'' or hurting the overall robustness of the engine? How is it possible to find the right balance between preventing health anxiety and the potential delay of important medical treatment when users are confronted with only benign explanations when they are actually sick?\n\nAlthough we studied the behavior many searchers, one limitation is that we used a single search engine, meaning our findings pertaining to user behavior might not generalize. Another limitation is the imperfect measurement of {\\it SERPConcern}. Due to the fact that some queries occur very frequently (e.g. ``headache''), any error incurred by our scoring function on the top results for those queries might have a big impact on the outcome. As mentioned earlier, another limitation is that we only used data on queries that were related to headache. Finally, query logs offer only a limited view on health concerns, and we need to work with searchers directly to understand the motivations behind their search behaviors.\n\n\\section{CONCLUSIONS}\n\nWe presented a log-based longitudinal study of medical search behavior on the web. We investigated the use of medical query formulation as a lens onto searchers health concerns and found that those features were predictive for this task. We evaluated how a major search engine responds to changes in medical query formulation and saw that there were some trends and a lot of variance, which might have adverse affects by way of misinformation. We showed that a significant tendency for medically concerned users to view concerning content makes it important for engines to manage this effect (e.g., by considering estimated level of health concern as a personalization feature). Finally, we raised the need for detailed study of the impact of SERP composition on users' future behavior. We believe that our results can function as an initial guide for developing practical tools for better health search as well as to inform deeper investigations into health concerns and anxieties on the Web and in general.\n\n\n\n\n\\bibliographystyle{abbrv}\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\n\\begin{sloppypar}\nMarkov chain Monte Carlo (MCMC) methods enable importance sampling from complicated, concentrated probability distributions.\nThe Metropolis-Hastings (MH) algorithm~\\cite{Metropolis1953Equation, Hastings1970Monte} is the prototype for a class of MCMC algorithms where transition proposals are accepted or rejected to generate each sample.\nIts applicability to distributions where probability ratios, but not absolute probabilities, are easily computed motivates its broad usage.\nSince the complexity and scale of these applications routinely push against computational resource limits, practical techniques for minimizing storage requirements and execution time are important aspects of implementations.\nIdeally, the full sequence of samples is recorded in a compact format that facilitates analysis and interpretation.\n\\end{sloppypar}\n\nCurrent implementations of Metropolis-Hastings-class algorithms typically reduce the information content of the full sample sequence in two ways before it is recorded.\nIn the first, a small fraction of samples are recorded in complete detail, but the rest are discarded.\nThis enables arbitrary post-analysis because all degrees of freedom are available for the retained samples.\nHowever, when sampling systems with many degrees of freedom, retaining even a small fraction of samples requires considerable storage capacity.\nEven if capacity is abundant, disk throughput would limit the fraction of samples that could be recorded without having disk write operations dominate the execution time.\nIn the second pattern, averages, moments, histograms, or other quantities are accumulated during simulation and recorded upon completion.\nThis is parsimonious with respect to demands on storage bandwidth and capacity, but decisions about which quantities to accumulate, their frequencies of accumulation, and the number of samples to discard due to starting configuration bias must be made in advance.\nAltering these choices requires repeating the entire simulation.\n\nThese tradeoffs and losses of information complicate the usage of MH-class algorithms in real applications.\nA typical experience from our lab provides an illustrative example in the area of biomolecular simulation.\nA study of conformational and dimerization equilibria~\\cite{Williamson2010Modulation} involved a set of $\\sim$$10$ proteins, each consisting of $\\sim$$10^3$ interacting atoms and modelled at $\\sim$$10$ different temperatures either individually or in pairs.\nAt least three independent replicate simulations were performed for each condition, with each replicate generating $\\sim$$10^8$ samples using MH.\nWhile some quantities of interest had their expectation values accumulated during the initial simulations, only 1 in $\\sim$$10^4$ samples had their full set of atomic coordinates recorded.\nDespite using Gromacs XTC compression with limited precision, these sparse recordings consumed a total of $\\sim$$10^{11}$ bytes of storage.\nSubsequent analyses could only be performed by reanalyzing these recorded samples or repeating the entire set of simulations while accumulating expectations for the new quantities of interest.\nGiven that the simulations consumed $\\sim$$10^5$ total hours of cpu time, both options were suboptimal in terms of efficient usage of computational resources.\n\nThese problems can be completely avoided.\nHere, we describe a simple method for recording the full sample sequence from a MH-class simulation that stores one bit per sample.\nThe method operates independently from any details of the system and the transition proposals, and is therefore generally applicable.\n\n\\section{Description of the method}\n\nRecording a bit string of 1's for acceptance and 0's for rejection, with one bit per transition, suffices to preserve the complete information content of all samples generated during one run of a MH-class algorithm.\nThis method follows naturally from an information theoretic perspective: since all influence from the underlying distribution is reduced to a binary accept-or-reject decision, MH-class algorithms can be viewed as communication channels with capacity of one bit per sample~\\cite[chapter 30.5]{MacKay2002Information}.\nAn existing implementation of a MH-class algorithm can be easily modified to perform this recording operation during each iteration.\nBuffered output is necessary because file systems do not allow writing of individual bits, and also helps to reduce the frequency of disk writes.\n\nRegenerating the full sequence of samples for subsequent reanalysis requires the recorded bit string, the starting state, and any pseudorandom number generator seed(s) from the original simulation.\nThe original simulation code should be reused to guarantee that the original sequence of transition proposals is recapitulated, but modified to use the recorded bit string for deciding whether to accept or reject proposals.\nSince the acceptance criterion no longer needs to be evaluated, all calculations involved in computing ratios of sample weights or proposal probabilities can be skipped during reanalysis.\nHowever, as with the original simulation run, all degrees of freedom for every sample are available for computing and accumulating distributions of any quantity of interest.\n\nCare must be taken to avoid corrupting the sequence of samples during reanalysis.\nThe stream of pseudorandom numbers generated during reanalysis must be identical to the original simulation's stream.\nIn particular, if a single pseudorandom number stream is used for generating transition proposals as well as evaluating the acceptance criterion, the pseudorandom number that would have been used for the acceptance criterion must be generated and discarded during each iteration.\nIf analysis routines themselves make use of pseudorandom numbers, they must generate their own independent streams.\nInsidious platform dependencies are another source of potential corruption when recordings generated on one computer are reanalyzed on another.\nFor example, the implementation of floating point arithmetic differs between processor architectures and compilers.\nThis source of error should be eliminated by adhering to best practices in the coding of floating point operations~\\cite{Goldberg1991What} or writing unit tests that assert platform-specific assumptions are valid during both original simulation and reanalysis runs.\n\n\\section{Demonstration of the method}\n\nTo create a nontrivial demonstration of this recording method, we implemented a MH simulation of a colloidal suspension of charged spherical particles.\nThis three-dimensional off-lattice system consists of $10^3$ particles confined within a spherical droplet.\nEach particle has two charge sites that freely diffuse on the particle's surface.\nThe potential energy is the sum of a pairwise Lennard-Jones potential between particle centers and a pairwise Debye-H\\\"{u}ckel screened electrostatic potential between charge sites.\nIntra-particle and inter-particle electrostatic interactions are screened using different dielectric constants.\nTransition proposals consist of local or full randomization of one particle's center coordinates or charge site positions.\nAn executable Java Archive file containing source code and compiled class files for this demonstration is available as supplementary material.\nTo run it on any computer where Java Runtime Environment version 6 is installed, execute \\texttt{java~-jar~RecordingDemonstration.jar} at a command line.\n\nA comparison with typical alternatives highlights the efficiency of the proposed recording method.\nA straightforward format would be a recording of all seven degrees of freedom $(x, y, z, \\theta_1, \\phi_1, \\theta_2, \\phi_2)$ for all particles.\nWhile this is obviously inefficient, it is simple and easy to parse, facilitating interoperability with other software.\nStoring the updated values, if any, for only the changed degrees of freedom during a transition would be significantly more efficient.\nHowever, this method would tie the recording format to the choice of transition proposals; modifying the simulation to use more sophisticated transitions that simultaneously perturb multiple particles would require redesigning the format.\nAn even more efficient method would record the potential energy (or more generally, the weight) of each sample.\nAs with the proposed method, this would require preserving the simulation code, but would enable full reconstruction of the original samples without calculating any weights.\nTable~\\ref{table:storagetable} compares the storage efficiency of these recording schemes to the proposed one.\n\\texttt{RecordingDemonstration.jar} can derive any of the other recordings starting from the bit string recording, proving that the bit string (along with simulation code and starting state) retains all information about the sample sequence.\n\n\\begin{table}[htbp]\n\\input{storagetable}\n\\label{table:storagetable}\n\\end{table}\n\n\\begin{sloppypar}\nSince the bit string recording obviates all sample weight and proposal probability ratio calculations, iterating over the sample sequence takes significantly less time than generating it.\nTo enable benchmarking, \\texttt{RecordingDemonstration.jar} provides two analyses for the demonstration system.\nThe first is a histogram of the central angle between intra-particle charge sites and the second is a pairwise distance histogram between particle centers.\nThe histograms are accumulated once every 100 and 1000 steps, respectively.\nTable~\\ref{table:timingtable} compares the execution time for performing both analyses during the original simulation versus using the bit string recording.\nNote that energy evaluations are efficiently implemented such that only changing terms are computed; each iteration after the first of the original simulation performs a number of computations that is linear, rather than quadratic, in the number of particles.\nHowever, once the bit string is recorded, energy evaluations are skipped; each iteration only needs to update the degrees of freedom for a single particle and therefore executes in constant time.\n\\end{sloppypar}\n\n\\begin{table}[htbp]\n\\input{timingtable}\n\\label{table:timingtable}\n\\end{table}\n\n\\section{Discussion}\n\nThis method enables efficient reanalysis as long as the time spent generating and executing transition proposals is small compared to the time spent computing ratios of sample weights and proposal probabilities.\nFortunately, this condition is naturally satisfied because sample weights are generally determined by an energy function or other interaction between degrees of freedom that is more expensive to update than the degrees of freedom themselves.\nBy recording the results of the most computationally expensive component of MH-class algorithm implementations, the method proposed here approaches the minimum achievable limits on both storage consumption and execution time.\n\nA closely related alternative method would be to record the full sequence of sample weights in addition to the accept \/ reject decision.\nIn some applications, the sample weights themselves are important subjects of analysis, and rederiving them from the samples during reanalysis would be computationally expensive.\nFor instance, in a simulation of a physical system in the canonical ensemble, one may wish to calculate the heat capacity and other properties of the energy distribution.\nAssuming that weights are represented as IEEE 754 double precision floating point numbers, as in Table~\\ref{table:storagetable}, each sample would require an additional 64 bits of storage.\nWhile much less compact than the bit string recording alone, it would still be far more efficient than most alternatives.\n\nNote that the size of the recording might be further reduced by compressing it using a lossless algorithm.\nOne bit per sample is an \\emph{upper} bound on the entropy rate~\\cite[chapter 30.5]{MacKay2002Information}; if the acceptance rate of the simulation is not exactly 1\/2, the recorded bit string will contain biases that a compression algorithm can exploit.\n\nThe compactness of this recording method makes it susceptible to corruption: during reanalysis, a single incorrectly recorded bit contaminates all subsequent samples.\nAn effective solution would be to encode the raw bit string using error correcting codes (ECC) that provide robustness at the cost of increased storage consumption~\\cite[chapter 11]{MacKay2002Information}.\nIn fact, employment of Reed-Solomon~\\cite{Reed1960Polynomial} and low-density parity-check (LDPC)~\\cite{Gallager1962Lowdensity} codes at the hardware level is already widespread among designers of random access memory, persistent storage devices, and networking interfaces.\nTherefore, ECC at the software level would constitute a redundant layer of protection against error.\nNon-redundant detection of error can be achieved by recording all degrees of freedom for a small number of checkpoint samples, and verifying equality for the corresponding samples generated during reanalysis.\nAlternatively, the common practice of comparing cryptographic hash function outputs can be applied to verifying the integrity of recorded bit strings.\n\nThis method is compatible with descendants of the MH algorithm that retain its general accept-or-reject architecture.\nExamples include expanded ensemble techniques~\\cite{Lyubartsev1992New}, replica exchange Monte Carlo~\\cite{Mitsutake2001Generalizedensemble}, multiple-try Metropolis~\\cite{Liu2000MultipleTry}, simulated annealing~\\cite{Kirkpatrick1983Optimization}, and simulated tempering~\\cite{Marinari1992Simulated}.\nFor some algorithms, such as Wang-Landau sampling~\\cite{Wang2001Efficient}, efficient post-analysis would require recording the system energy at every iteration as described above.\n\n\\section{Conclusion}\n\nThe method described here is a simple and effective approach to data storage and representation in the design of MH-class algorithm implementations.\nIt facilitates decoupling of simulation design from the constraints of data analysis.\nThe information theoretic perspective through which this method was conceived deserves broader appreciation in the development of Monte Carlo algorithms.\n\n\\section{Acknowledgements}\n\nThis work was supported by National Science Foundation MCB 0718924 and National Institutes of Health \u2013 National Institute of General Medical Sciences 5T32GM008802.\n\n\n\n\n\n\n\\bibliographystyle{model1-num-names}\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{INTRODUCTION}\n\nObservations of high-energy radiation from a variety of astrophysical\nsources imply the presence of significant populations of nonthermal\n(often relativistic) particles. Much of the observational data is\ncharacterized by strong variability on very short timescales. Nonthermal\ndistributions are naturally produced via the Fermi process, in which\nparticles interact with scattering centers moving systematically and\/or\nrandomly. The first-order Fermi mechanism treats particle acceleration\nin converging flows, such as shocks, and is thought to be important at\nthe Earth's bow shock, in certain classes of solar flares, for cosmic-ray\nacceleration by supernova remnants, and in sources with relativistic\noutflows, including blazars and $\\gamma$-ray bursts \\citep[for reviews,\nsee][]{be87,kir94}. The second-order process, as it was originally\nconceived by Fermi, involved the stochastic acceleration of particles\nscattering with randomly moving magnetized clouds \\citep{fer49}. Later\nrefinements of this idea replaced the magnetized clouds with\nmagnetohydrodynamic (MHD) waves \\citep[e.g.,][]{mel80}. The\nsecond-order, stochastic Fermi process now finds application in a wide\nrange of astrophysical settings, including solar flares\n\\citep{mgr90,lpm04a}, clusters of galaxies \\citep{sst87,pet01,bru04},\nthe Galactic center \\citep{lpm04b,ad04}, and $\\gamma$-ray bursts\n\\citep{wax95,dh01}.\n\nThe standard approach to modeling the acceleration of nonthermal\nparticles via interactions with MHD waves involves obtaining the\nsolution to a steady-state transport equation that incorporates\ntreatments of systematic and stochastic Fermi acceleration, radiative\nlosses, and particle escape \\citep[e.g.,][]{sch84,ss89,lmp06}. However,\nthe prevalence of variability on short timescales in many sources calls\ninto question the validity of the steady-state interpretation.\n\\citet{pp95} provided a comprehensive review of the various\ntime-dependent solutions that have been derived in the past 30 years. In\nmost cases, it is assumed that the momentum diffusion coefficient,\n$D(p)$, and the mean particle escape timescale, $t_{\\rm esc}(p)$, have\npower-law dependences on the particle momentum $p$. Although the set of\nexisting solutions covers a broad range of possible values for the\nassociated power-law indices, there are certain physically interesting\ncases for which no analytical solution is currently available. For\nexample, \\citet{dml96} have shown that for the stochastic acceleration\nof relativistic particles due to resonant interactions with plasma waves\nin black-hole magnetospheres, one obtains $D(p) \\propto p^q$ and $t_{\\rm\nesc}(p) \\propto p^{q-2}$, where $q$ is the index of the wavenumber\ndistribution (see \\S~2). The analytical solution for the time-dependent\nGreen's function in this situation is of particular physical interest,\nbut it has not appeared in the previous literature. This has motivated\nus to reexamine the associated Fokker-Planck transport equation for this\ncase, and to obtain a new family of closed-form solutions for the\nsecular Green's function. The resulting expression, describing the\nevolution of an initially monoenergetic particle distribution,\ncomplements the set of solutions discussed by \\citet{pp95}. Our primary\ngoal in this paper is to present a detailed derivation of the exact\nanalytical solution and to demonstrate some of its key properties.\nDetailed applications of our results to the modeling of particle\nacceleration in black-hole accretion coronae and other astrophysical\nenvironments will be presented in a separate paper.\n\nThe remainder of the paper is organized as follows. In \\S~2 we review\nthe fundamental equations governing the acceleration of charged\nparticles interacting with plasma waves. The transport equation is\nsolved in \\S~3 to obtain the time-dependent Green's function, and\nillustrative results are presented in \\S~4. The astrophysical implications\nof our work are summarized in \\S~5, and additional mathematical details are\nprovided in the Appendix.\n\n\\section{FUNDAMENTAL EQUATIONS}\n\nCharged particles in turbulent astrophysical plasmas are expected to be\naccelerated via interactions with whistler, fast-mode, and Alfv\\'en\nwaves propagating in the magnetized gas. Here we consider a simplified\nisotropic description of the wave energy distribution, denoted by\n$W(k)$, where $W(k) dk$ represents the energy density due to waves with\nwavenumber between $k$ and $k+dk$. The transport formalism we consider\nassumes a power-law distribution for the wave-turbulence spectrum, which\nimplies definite relations between the momentum diffusion coefficient\nand the momentum-dependent escape timescale\n\\citep[e.g.,][]{mel74,sch89a,sch89b,dml96}. MHD waves injected over a\nnarrow range of wavenumber cascade to larger wavenumbers, forming a\nKolmogorov or Kraichnan power spectrum over the inertial range with\n$W(k) \\propto k^{-q}$, where $q=5\/3$ and $q=3\/2$ for the Kolmogorov and\nKraichnan cases, respectively \\citep[e.g.,][]{zm90}. The specific forms\nwe will adopt for the transport coefficients in \\S~2.2 are based on the\nphysics of the resonant scattering processes governing the wave-particle\ninteractions \\citep[see][]{dml96}. We assume that the nonthermal\nparticle distribution is isotropic, and focus on a detailed treatment of\nthe propagation of particles in momentum space due to wave-particle\ninteractions. The spatial aspects of the transport (i.e., the\nconfinement of the particles in the acceleration region) will be treated\nin an approximate manner using a momentum-dependent escape term.\n\n\\subsection{Transport Equation}\n\nThe fundamental transport equation describing the propagation of\nparticles in the momentum space can be written in the flux-conservation\nform \\citep[e.g.,][]{bec92,sch89a,sch89b}\n\\begin{equation}\n{\\partial f \\over \\partial t}\n= - {1 \\over p^2} {\\partial \\over \\partial p}\\left\\{\np^2 \\left[\nA(p) \\, f - D(p) {\\partial f \\over \\partial p}\\right]\n\\right\\}\n- {f \\over t_{\\rm esc}(p)}\n+ {S(p,t) \\over 4 \\pi p^2}\n\\ ,\n\\label{eq1}\n\\end{equation}\nwhere $p$ is the particle momentum, $f(p,t)$ is the particle\ndistribution function, $D(p)$ denotes the momentum diffusion\ncoefficient, $t_{\\rm esc}(p)$ is the mean escape time, $S(p,t)$\nrepresents particle sources, and $A(p)$ describes any additional,\nsystematic acceleration or loss processes, such as shock acceleration or\nsynchrotron\/inverse-Compton emission. The quantity in square brackets in\nequation~(\\ref{eq1}) describes the flux of particles through the\nmomentum space \\citep{tnj71}, and the source term is defined so that\n$S(p,t) \\, dp \\, dt$ gives the number of particles injected into the\nplasma per unit volume between times $t$ and $t + dt$ with momenta\nbetween $p$ and $p + dp$. The total particle number density $n(t)$ and\nenergy density $U(t)$ of the distribution $f(p,t)$ are computed using\n\\begin{equation}\nn(t) = \\int_0^\\infty 4 \\pi \\, p^2 \\, f(p,t) \\, dp\n\\ , \\ \\ \\ \\ \\ \nU(t) = \\int_0^\\infty 4 \\pi \\, \\epsilon \\, p^2 \\, f(p,t)\n\\, dp \\ ,\n\\label{eq2}\n\\end{equation}\nwhere the particle kinetic energy $\\epsilon$ is related to the Lorentz factor\n$\\gamma$ and the particle momentum $p$ by\n\\begin{equation}\n\\epsilon = (\\gamma-1) \\, m c^2 \\ , \\ \\ \\ \\ \\ \n\\gamma=\\left({p^2 \\over m^2 c^2} + 1\\right)^{1\/2} \\ ,\n\\label{eq3}\n\\end{equation}\nand $m$ and $c$ denote the particle rest mass and the speed of light,\nrespectively.\n\nRather than solve equation~(\\ref{eq1}) directly to determine $f(p,t)$\nfor a given source term $S(p,t)$, it is more instructive to first solve\nfor the Green's function, $f_{_{\\rm G}}(p_0,p,t_0,t)$, which satisfies the\nequation\n\\begin{equation}\n{\\partial f_{_{\\rm G}} \\over \\partial t}\n= - {1 \\over p^2} {\\partial \\over \\partial p}\\left\\{\np^2 \\left[A(p) \\, f_{_{\\rm G}}\n- D(p) {\\partial f_{_{\\rm G}} \\over \\partial p}\n\\right]\\right\\}\n- {f_{_{\\rm G}} \\over t_{\\rm esc}(p)}\n+ {\\delta(p-p_0) \\, \\delta(t-t_0) \\over 4 \\pi p_0^2}\n\\ ,\n\\label{eq4}\n\\end{equation}\nwhere $p_0$ is the initial momentum and $t_0$ is the initial time.\nThe source term in this equation corresponds to the injection of a single\nparticle per unit volume with momentum $p_0$ at time $t_0$. The particle\nnumber and energy densities associated with the Green's function are\ngiven by\n\\begin{equation}\nn_{_{\\rm G}}(t) = \\int_0^\\infty 4 \\pi \\, p^2 \\,\nf_{_{\\rm G}}(p_0,p,t_0,t) \\, dp\n\\ , \\ \\ \\ \\ \\ \nU_{_{\\rm G}}(t) = \\int_0^\\infty 4 \\pi \\, \\epsilon \\, p^2 \\,\nf_{_{\\rm G}}(p_0,p,t_0,t) \\, dp\n\\ .\n\\label{eq5}\n\\end{equation}\nOnce the Green's function solution has been determined, the {\\it\nparticular solution} associated with an arbitrary source distribution\n$S(p,t)$ can be computed using the integral convolution\n\\citep[e.g.,][]{bec03}\n\\begin{equation}\nf(p,t) = \\int_0^t \\int_0^\\infty\nf_{_{\\rm G}}(p_0,p,t_0,t) \\,\nS(p_0,t_0) \\, dp_0 \\, dt_0\n\\ ,\n\\label{eq6}\n\\end{equation}\nwhere we have assumed that the particle injection begins at time $t=0$\nand no particles are present in the plasma for $t < 0$.\n\n\\subsection{Transport Coefficients}\n\nIn Appendix~A.1, we demonstrate that for\narbitrary particle energies, the mean rate of change of the\nparticle momentum due to stochastic acceleration is related to the\nmomentum diffusion coefficient $D(p)$ via\n\\begin{equation}\n\\Big\\langle {dp \\over dt} \\Big\\rangle \\Big|_{\\rm stoch} = {1 \\over p^2}\n{d \\over dp}\\left(p^2 D\\right)\n\\ .\n\\label{eq8}\n\\end{equation}\nThe corresponding result for the mean rate of change of the kinetic energy\nis \\citep[see Appendix~A.2 and][]{mr89}\n\\begin{equation}\n\\Big\\langle {d\\epsilon \\over dt} \\Big\\rangle \\Big|_{\\rm stoch}\n= {1 \\over p^2} {d \\over dp} \\left(p^2 v D \\right) \\ ,\n\\label{eq8b}\n\\end{equation}\nwhere $v$ is the particle speed. If the MHD wave spectrum has the\npower-law form $W \\propto k^{-q}$ associated with Alfv\\'en and fast-mode\nwaves, then the momentum dependences of the diffusion coefficient $D(p)$\nand the mean escape time $t_{\\rm esc}(p)$ describing the resonant\npitch-angle scattering of relativistic particles are given by\n\\citep[e.g.,][]{dml96,mr89}\n\\begin{equation}\nD(p) = D_* \\, m^2 c^2 \\left({p \\over m c}\n\\right)^q \\ , \\ \\ \\ \\ \\ \\ \\\nt_{\\rm esc}(p) = t_* \\left({p \\over m c}\n\\right)^{q-2} \\ ,\n\\label{eq7}\n\\end{equation}\nwhere $D_* \\propto {\\rm s}^{-1}$ and $t_* \\propto {\\rm s}$ are\nconstants. We shall focus on transport scenarios with $q \\le 2$, so that\n$t_{\\rm esc}$ is either a decreasing or constant function of the\nmomentum $p$. In order to maintain the physical validity of the\nescape-timescale approach used here, we must require that $t_{\\rm esc}$\nexceed the light-crossing time $L\/c$ for a source with size $L$. This\nimplies a fundamental upper limit to the particle momentum when $q < 2$.\n\nBy combining equations~(\\ref{eq8}) and (\\ref{eq7}), we find that the\nmean rate of change of the momentum for relativistic particles\naccelerated stochastically by MHD waves is given by\n\\begin{equation}\n\\Big\\langle {dp \\over dt} \\Big\\rangle\\Big|_{\\rm stoch} = (q+2) \\, D_*\n\\, m c \\left({p \\over m c}\\right)^{q-1}\n\\ .\n\\label{eq9}\n\\end{equation}\nFor simplicity, we shall assume that the momentum dependence of the\nadditional, systematic loss\/acceleration processes appearing in the transport\nequation (\\ref{eq1}), described by the coefficient $A(p)$, mimics that of the\nstochastic acceleration (eq.~[\\ref{eq9}]). We therefore write\n\\begin{equation}\nA(p) = A_* \\, m c \\left({p \\over m c}\\right)^{q-1}\n\\ ,\n\\label{eq10}\n\\end{equation}\nwhere the constant $A_* \\propto {\\rm s}^{-1}$ determines the positive\n(negative) rate of systematic acceleration (losses). Note that this\nformulation precludes the treatment of loss processes with a quadratic\nenergy dependence (e.g., inverse-Compton or synchrotron) since that\nwould imply $q=3$, which is outside the range considered here. However,\nfirst-order Fermi acceleration at a shock or energy losses due to\nCoulomb collisions can be treated by setting $q=2$ with $A_*$ either\npositive or negative, respectively. This suggests that the results\nobtained here are relevant primarily for the transport of energetic\nions. However, even in this application one needs to bear in mind that\nsynchrotron and inverse-Compton losses will become dominant at\nsufficiently high energies \\citep[e.g.,][]{sch84,ss89,lmp06}.\n\nIt is convenient to transform to the dimensionless momentum and time variables\n$x$ and $y$, defined by\n\\begin{equation}\nx \\equiv {p \\over m c} \\ , \\ \\ \\ \\ \\ \\ \\ \ny \\equiv D_* \\, t\n\\ ,\n\\label{eq11}\n\\end{equation}\nin which case the transport equation~(\\ref{eq4}) for the Green's function becomes\n\\begin{equation}\n{\\partial f_{_{\\rm G}} \\over \\partial y}\n= {1 \\over x^2} {\\partial \\over \\partial x}\\left(\nx^{2+q} \\, {\\partial f_{_{\\rm G}} \\over \\partial x}\\right)\n- {a \\over x^2} {\\partial \\over \\partial x} \\left(\nx^{1+q} \\, f_{_{\\rm G}}\\right) - \\theta \\, x^{2-q} \\, f_{_{\\rm G}}\n+ {\\delta(x-x_0) \\, \\delta(y-y_0) \\over 4 \\pi m^3 c^3 x_0^2}\n\\ ,\n\\label{eq12}\n\\end{equation}\nwhere\n\\begin{equation}\nx_0 \\equiv {p_0 \\over m c} \\ , \\ \\ \\ \\ \\ \\ \\ \ny_0 \\equiv D_* \\, t_0\n\\ ,\n\\label{eq13}\n\\end{equation}\nand we have introduced the dimensionless constants\n\\begin{equation}\na \\equiv {A_* \\over D_*} \\ , \\ \\ \\ \\ \\ \\ \\ \n\\theta \\equiv {1 \\over D_* \\, t_*}\n\\ .\n\\label{eq14}\n\\end{equation}\nNote that $x$ equals the particle Lorentz factor in applications involving\nultrarelativistic particles. The constant $a$ describes the relative\nimportance of systematic gains or losses compared with the stochastic\nprocess.\n\n\\subsection{Fokker-Planck Equation}\n\nAdditional physical insight can be obtained by reorganizing equation~(\\ref{eq12})\nin the Fokker-Planck form\n\\begin{equation}\n{\\partial N_{_{\\rm G}} \\over \\partial y}\n= {\\partial^2 \\over \\partial x^2} \\left(\nx^q \\, N_{_{\\rm G}}\\right)\n- {\\partial \\over \\partial x} \\left[(q+2+a) \\, x^{q-1}\n\\, N_{_{\\rm G}}\\right] - \\theta \\, x^{2-q} \\, N_{_{\\rm G}}\n+ {\\delta(x-x_0) \\, \\delta(y-y_0)}\n\\ ,\n\\label{eq15}\n\\end{equation}\nwhere we have defined the Green's function number distribution $N_{_{\\rm G}}$\nusing\n\\begin{equation}\nN_{_{\\rm G}}(x_0,x,y_0,y) \\equiv 4 \\pi \\, m^3 c^3 \\, x^2 \\, f_{_{\\rm G}}(x_0,x,y_0,y)\n\\ .\n\\label{eq16}\n\\end{equation}\nThe Fokker-Planck coefficients appearing in equation~(\\ref{eq15}), which\ndescribe the evolution of the particle distribution due to the influence\nof stochastic and systematic processes, are given by \\citep{rei65}\n\\begin{equation}\n{1 \\over 2}{d\\sigma^2 \\over dy} = x^q \\ , \\ \\ \\ \\ \\ \n\\Big\\langle{dx\\over dy}\\Big\\rangle = (q+2+a) \\, x^{q-1} \\ ,\n\\label{eq17}\n\\end{equation}\nwhere the first coefficient describes the ``broadening'' of the distribution\ndue to momentum space diffusion, and the second represents the mean ``drift''\nof the particles (i.e., the average acceleration rate).\n\nEquation~(\\ref{eq15}) is equivalent to equation~(49) from \\citet{pp95}\nif we set their parameters $a_{\\rm pp} = 2+a$ and $s_{\\rm pp} = q-2$,\nwhere $s_{\\rm pp}$ denotes the power-law index describing the momentum\ndependence of the escape timescale. Our particular choice for $s_{\\rm\npp}$ reflects the physics of the resonant wave-particle interactions, as\nrepresented by equations~(\\ref{eq7}). The Fokker-Planck form of\nequation~(\\ref{eq15}) clearly reveals the fundamental nature of the\ntransport process. In particular, we note that in the limit $a \\to 0$,\nthe drift coefficient $\\langle dx\/dy\\rangle$ reduces to the purely\nstochastic result (eq.~[\\ref{eq9}]), as expected when systematic\ngains\/losses are excluded from the problem. The total number and energy\ndensities are computed in terms of $N_{_{\\rm G}}$ using (cf. eq.~[\\ref{eq5}])\n\\begin{equation}\nn_{_{\\rm G}}(y) = \\int_0^\\infty N_{_{\\rm G}}(x_0,x,y_0,y) \\, dx\n\\ , \\ \\ \\ \\ \\ \nU_{_{\\rm G}}(y) = \\int_0^\\infty \\epsilon \\, x \\, N_{_{\\rm G}}(x_0,x,y_0,y) \\, dx\n\\ ,\n\\label{eq18}\n\\end{equation}\nwhere (see eq.~[\\ref{eq3}])\n\\begin{equation}\n\\epsilon = m c^2 \\left[\\left(x^2 + 1\\right)^{1\/2}-1\\right]\n\\ .\n\\label{eq19}\n\\end{equation}\nSince the physical situation considered here corresponds to the injection\nof a single particle per unit volume at ``time'' $y=y_0$, it follows that\n$n_{_{\\rm G}}(y_0)=1$.\n\n\\section{SOLUTION FOR THE GREEN'S FUNCTION}\n\nThe result obtained for the Green's function number distribution\n$N_{_{\\rm G}}$ has two important special cases depending on whether $q = 2$\nor $q < 2$. \\citet{pp95} obtained the exact solution to\nequation~(\\ref{eq15}) for the hard sphere case \\citep[][]{ram79} with\n$q=2$ and their parameter $s_{\\rm pp}=0$, corresponding to a\nmomentum-independent escape timescale. In this section we derive the\nexact solution to the time-dependent problem with $q < 2$ and $s_{\\rm\npp}=q-2$, which describes the physics of the wave-particle interactions\n(see eqs.~[\\ref{eq7}]).\n\n\\subsection{Laplace Transformation}\n\nWe define the Laplace transformation of $N_{_{\\rm G}}$ using\n\\begin{equation}\nL(x_0,x,s) \\equiv \\int_0^\\infty e^{-s y} \\, N_{_{\\rm G}}(x_0,x,y_0,y) \\, dy\n\\ .\n\\label{eq20}\n\\end{equation}\nBy operating on the Fokker-Planck equation~(\\ref{eq15}) with\n$\\int_0^\\infty e^{-s y} \\, dy$, we obtain\n\\begin{equation}\n{d^2 \\over d x^2} \\left(x^q \\, L \\right)\n- {d \\over d x} \\left[(q+2+a) \\, x^{q-1}\n\\, L \\right] - \\theta \\, x^{2-q} \\, L - s L\n= - e^{-s y_0} \\, \\delta(x-x_0)\n\\ ,\n\\label{eq21}\n\\end{equation}\nor, equivalently,\n\\begin{equation}\n{d^2 L \\over d x^2} + \\left({q - 2 - a \\over x} \\right)\n{d L \\over d x}\n+ \\left[{(1-q)(2+a) \\over x^2} - {\\theta \\over x^{2 q-2}}\n- {s \\over x^q} \\right] \\, L = - {e^{-s y_0} \\, \\delta(x-x_0) \\over x^q}\n\\ .\n\\label{eq22}\n\\end{equation}\nThis equation can be transformed into standard form by introducing the\nnew momentum variables\n\\begin{equation}\nz(x) \\equiv {2 \\sqrt{\\theta} \\over 2 - q} \\ x^{2-q} \\ , \\ \\ \\ \\ \\ \nz_0(x_0) \\equiv {2 \\sqrt{\\theta} \\over 2 - q} \\ x_0^{2-q} \\ .\n\\label{eq23}\n\\end{equation}\nAfter some algebra, we find that the equation for $L$ now becomes\n\\begin{equation}\nz^2 \\, {d^2 L \\over d z^2} + {a+1 \\over q-2} \\, z \\,\n{d L \\over d z} + \\left[{(1-q)(2+a) \\over (2-q)^2}\n- {z^2 \\over 4} - {s z \\over c_0 (2-q)^2}\\right] \\, L\n= - {c_0 \\, e^{-s y_0} \\, \\delta(z-z_0) \\over 2-q} \\, \\left({z \\over c_0}\n\\right)^{(3-2q)\/(2-q)}\n\\ ,\n\\label{eq24}\n\\end{equation}\nwhere\n\\begin{equation}\nc_0 \\equiv {2 \\sqrt{\\theta} \\over 2-q} \\ .\n\\label{eq25}\n\\end{equation}\nThe solutions to equation~(\\ref{eq24}) obtained for $z \\ne z_0$ that satisfy\nthe high- and low-energy boundary conditions are given by\n\\begin{equation}\nL(z_0,z,s) = e^{-z\/2} \\, z^{(a+2)\/(2-q)} \\,\n\\cases{\nA \\, U(\\alpha,\\beta,z) \\ , & $z \\ge z_0$ \\ , \\cr\nB \\, M(\\alpha,\\beta,z) \\ , & $z \\le z_0$ \\ , \\cr\n}\n\\label{eq26}\n\\end{equation}\nwhere $M$ and $U$ denote the confluent hypergeometric functions\n\\citep{as70}, and\n\\begin{equation}\n\\alpha \\equiv {s + (a+3) \\sqrt{\\theta} \\over 2 (2-q) \\sqrt{\\theta}}\n\\ , \\ \\ \\ \\ \\ \\ \\\n\\beta \\equiv {a+3 \\over 2-q} \\ .\n\\label{eq27}\n\\end{equation}\nThe constants $A$ and $B$ appearing in equation~(\\ref{eq26}) are determined\nby ensuring that the function $L$ is continuous at $z=z_0$, and that it also\nsatisfies the derivative jump condition\n\\begin{equation}\n\\lim_{\\varepsilon \\to 0} {dL \\over dz} \\bigg|_{z_0-\\varepsilon}\n^{z_0+\\varepsilon} = - {c_0 \\, e^{-s y_0} \\over (2-q) \\, z_0^2} \\,\n\\left({z_0 \\over c_0} \\right)^{(3-2q)\/(2-q)}\n= {- e^{-s y_0} \\over 2 \\, x_0 \\sqrt{\\theta}}\n\\ ,\n\\label{eq28}\n\\end{equation}\nobtained by integrating the transport equation~(\\ref{eq24}) with respect\nto $z$ in a small range around the source momentum $z_0$.\n\nThe constant $B$ can be eliminated by combining the continuity and derivative\njump conditions. After some algebra, the solution obtained for $A$ is\n\\begin{equation}\nA = - \\, {e^{-s y_0} \\, e^{z_0\/2} \\, z_0^{(a+2)\/(q-2)} \\, M(\\alpha,\\beta,z_0)\n\\over 2 \\, x_0 \\sqrt{\\theta} \\, W(z_0)}\n\\ ,\n\\label{eq29}\n\\end{equation}\nwhere $W(z)$ denotes the Wronskian, defined by\n\\begin{equation}\nW(z) \\equiv\nM(\\alpha,\\beta,z) {d \\over dz} U(\\alpha,\\beta,z)\n- U(\\alpha,\\beta,z) {d \\over dz} M(\\alpha,\\beta,z)\n\\ .\n\\label{eq30}\n\\end{equation}\nUsing equation~(13.1.22) from \\citet{as70}, we find that $W$ is given\nby the exact expression\n\\begin{equation}\nW(z) = - \\, {\\Gamma(\\beta) \\, z^{-\\beta} \\, e^z \\over\n\\Gamma(\\alpha)}\n\\ ,\n\\label{eq31}\n\\end{equation}\nwhich can be combined with equation~(\\ref{eq29}) and the continuity\nrelation to obtain\n\\begin{equation}\nA = {\\Gamma(\\alpha) \\, z_0^{\\beta+(a+2)\/(q-2)} \\, e^{-s y_0} \\, e^{-z_0\/2}\n\\, M(\\alpha,\\beta,z_0) \\over 2 \\, x_0 \\sqrt{\\theta} \\, \\Gamma(\\beta)}\n\\ ,\n\\label{eq32}\n\\end{equation}\n\\begin{equation}\nB = {\\Gamma(\\alpha) \\, z_0^{\\beta+(a+2)\/(q-2)} \\, e^{-s y_0} \\, e^{-z_0\/2}\n\\, U(\\alpha,\\beta,z_0) \\over 2 \\, x_0 \\sqrt{\\theta} \\, \\Gamma(\\beta)}\n\\ .\n\\label{eq33}\n\\end{equation}\nThe final solution for the Laplace transformation $L$ can therefore\nbe written as\n\\begin{equation}\nL(z_0,z,s) = {\\Gamma(\\alpha) \\, z_0^\\beta \\over \\Gamma(\\beta) \\,\n2 \\, x_0 \\sqrt{\\theta}} \\, \\left({z \\over z_0} \\right)^{(a+2)\/(2-q)}\n\\, e^{-s y_0} \\, e^{-(z+z_0)\/2} \\,\nM(\\alpha,\\beta,z_{\\rm min}) \\, U(\\alpha,\\beta,z_{\\rm max})\n\\ ,\n\\label{eq34}\n\\end{equation}\nwhere\n\\begin{equation}\nz_{\\rm min} \\equiv \\min(z,z_0) \\ , \\ \\ \\ \\ \\ \\ \\\nz_{\\rm max} \\equiv \\max(z,z_0)\n\\ .\n\\label{eq35}\n\\end{equation}\n\n\\subsection{Inverse Transformation}\n\n\\begin{figure}[t]\n\\hspace{35mm}\n\\includegraphics[width=85mm]{f1.eps}\n\\caption{Integration contour $C$ used to evaluate equation~(\\ref{eq36}).}\n\\label{fig1}\n\\end{figure}\n\nThe solution for the Green's function $N_{_{\\rm G}}$ can be found using the\ncomplex Mellin inversion integral, which states that\n\\begin{equation}\nN_{_{\\rm G}}(z_0,z,y_0,y) = {1 \\over 2 \\pi i} \\int_{\\gamma-i \\infty}\n^{\\gamma + i \\infty} e^{s y} \\, L(z_0,z,s) \\, ds\n\\ ,\n\\label{eq36}\n\\end{equation}\nwhere $\\gamma$ is chosen so that the line ${\\rm Re} \\, s = \\gamma$ lies\nto the right of any singularities in the integrand. The singularities\nare simple poles located where $|\\Gamma(\\alpha)| \\to \\infty$, which\ncorresponds to $\\alpha = -n$, with $n = 0,1,2,\\ldots$\nEquation~(\\ref{eq27}) therefore implies that there are an infinite\nnumber of poles situated along the real axis at $s = s_n$, where\n\\begin{equation}\ns_n = \\left[2 \\, (q-2) \\, n - a - 3\\right] \\sqrt{\\theta}\n\\ , \\ \\ \\ n = 0,1,2,\\ldots\n\\label{eq37}\n\\end{equation}\nWe can therefore employ the residue theorem to write\n\\begin{equation}\n\\oint e^{s y} \\, L(z_0,z,s) \\, ds\n= 2 \\pi i \\sum_{n=0}^\\infty \\ {\\rm Res}(s_n)\n\\ ,\n\\label{eq38}\n\\end{equation}\nwhere $C$ is the closed integration contour indicated in Figure~1 and\n${\\rm Res}(s_n)$ denotes the residue associated with the pole at $s=s_n$.\nBased on asymptotic analysis, we conclude that the contribution to the integral\ndue to the curved portion of the contour vanishes in the limit $R \\to \\infty$,\nand consequently we can combine equations~(\\ref{eq36}) and (\\ref{eq38}) to obtain\n\\begin{equation}\nN_{_{\\rm G}}(z_0,z,y_0,y) = \\sum_{n=0}^\\infty \\ {\\rm Res}(s_n)\n\\ .\n\\label{eq39}\n\\end{equation}\nHence we need only evaluate the residues in order to determine the solution\nfor the Green's function.\n\n\\subsection{Evaluation of the Residues}\n\nThe residues associated with the simple poles at $s=s_n$ are evaluated\nusing the formula \\citep[e.g.,][]{but68}\n\\begin{equation}\n{\\rm Res}(s_n) = \\lim_{s \\to s_n}\n(s-s_n) \\ e^{s y} \\, L(z_0,z,s)\n\\ .\n\\label{eq40}\n\\end{equation}\nSince the poles are associated with the function $\\Gamma(\\alpha)$\nin equation~(\\ref{eq34}), we will need to make use of the identity\n\\begin{equation}\n\\lim_{s \\to s_n} (s-s_n) \\ \\Gamma(\\alpha)\n= {(-1)^n \\over n!} \\, 2 \\, (2-q) \\, \\sqrt{\\theta}\n\\ ,\n\\label{eq41}\n\\end{equation}\nwhich follows from equations~(\\ref{eq27}) and (\\ref{eq37}). Combining\nequations~(\\ref{eq34}), (\\ref{eq40}), and (\\ref{eq41}), we find that the\nresidues are given by\n\\begin{equation}\n{\\rm Res}(s_n) = {(-1)^n \\, (2-q) \\, e^{s_n (y-y_0)} \\, z_0^\\beta \\over\nn! \\, \\Gamma(\\beta) \\, x_0} \\, \\left({z \\over z_0}\\right)^{(a+2)\/(2-q)}\n\\, e^{-(z+z_0)\/2} \\, M(-n,\\beta,z_{\\rm min}) \\, U(-n,\\beta,z_{\\rm max})\n\\ .\n\\label{eq42}\n\\end{equation}\nBased on equations~(13.6.9) and (13.6.27) from Abramowitz \\& Stegun\n(1970), we note that the confluent hypergeometric functions appearing in\nthis expression reduce to Laguerre polynomials, and therefore our result\nfor the residue can be rewritten after some simplification as\n\\begin{equation}\n{\\rm Res}(s_n) = {n! \\, e^{s_n (y-y_0)} \\, z_0^\\beta \\, (2-q)\n\\over \\Gamma(\\beta+n) \\, x_0} \\, \\left({z \\over z_0}\\right)^{(a+2)\/(2-q)}\n\\, e^{-(z+z_0)\/2} \\, P_n^{(\\beta-1)}(z) \\, P_n^{(\\beta-1)}(z_0)\n\\ ,\n\\label{eq43}\n\\end{equation}\nwhere $P_n^{(\\beta-1)}(z)$ denotes the Laguerre polynomial.\n\n\\subsection{Closed-Form Expression for the Green's Function}\n\nThe result for the Green's function number distribution $N_{_{\\rm G}}$ is obtained\nby summing the residues (see eq.~[\\ref{eq39}]), which yields\n\\begin{equation}\nN_{_{\\rm G}}(z_0,z,y_0,y) = \\sum_{n=0}^\\infty\n\\ {n! \\, e^{s_n (y-y_0)} \\, z_0^\\beta \\, (2-q)\n\\over \\Gamma(\\beta+n) \\, x_0} \\, \\left({z \\over z_0}\\right)^{(a+2)\/(2-q)}\n\\, e^{-(z+z_0)\/2} \\, P_n^{(\\beta-1)}(z) \\, P_n^{(\\beta-1)}(z_0)\n\\ .\n\\label{eq44}\n\\end{equation}\nThis convergent sum is a useful expression for the Green's function.\nHowever, further progress can be made by employing the bilinear\ngenerating function for the Laguerre polynomials, given by\nequation~(8.976) from \\citet{gr80}. After some algebra, we find that the\ngeneral closed-form solution can be written in the form\n\\begin{equation}\nN_{_{\\rm G}}(x_0,x,y_0,y) = {2-q \\over x_0} \\left({x \\over x_0}\\right)^{(a+1)\/2}\n\\, {\\sqrt{z z_0 \\xi} \\over 1 - \\xi} \\ \\exp\\left[-{(z + z_0)(1+\\xi) \\over 2\n\\, (1-\\xi)} \\right] \\, I_{\\beta-1}\\left({2 \\sqrt{z z_0 \\xi} \\over 1 - \\xi}\\right)\n\\ ,\n\\label{eq45}\n\\end{equation}\nwhere\n\\begin{equation}\nz(x) \\equiv {2 \\sqrt{\\theta} \\over 2 - q} \\ x^{2-q} \\ , \\ \\ \\ \\ \\ \nz_0(x_0) \\equiv {2 \\sqrt{\\theta} \\over 2 - q} \\ x_0^{2-q} \\ , \\ \\ \\ \\ \\ \n\\beta \\equiv {a+3 \\over 2 - q} \\ , \\ \\ \\ \\ \\ \n\\xi(y,y_0) \\equiv e^{2(q-2)(y-y_0)\\sqrt{\\theta}}\n\\ .\n\\label{eq46}\n\\end{equation}\nNote that the solution for $N_{_{\\rm G}}$ depends on the time parameters $y$\nand $y_0$ only through the ``age'' of the injected particles, $y-y_0$, as\nindicated by the form for $\\xi$. In the limit $\\theta \\to 0$,\ncorresponding to infinite escape time, the Green's function number\ndistribution reduces to\n\\begin{equation}\nN_{_{\\rm G}}(x_0,x,y_0,y) \\Big|_{\\theta=0} \\! = {(x x_0)^{(3-q)\/2} \\over (2-q)\n\\, x_0^2 \\, (y-y_0)}\n\\left(x \\over x_0 \\right)^{a\/2} \\!\\! \\exp\\left[-{(x^{2-q} + x_0^{2-q}) \\over\n(2-q)^2 (y-y_0)}\\right]\nI_{\\beta-1}\\!\\left[{2 (x x_0)^{(2-q)\/2} \\over (2-q)^2 (y-y_0)}\\right]\n\\ .\n\\label{eq47}\n\\end{equation}\n\nThe exact solution for the time-dependent Green's function describing the\nevolution of a monoenergetic initial spectrum with $q < 2$ given\nby equation~(\\ref{eq45}) represents an interesting new contribution to\nparticle transport theory. The corresponding solution for the hard-sphere\ncase with $q = 2$, given by equation~(43) of \\citet{pp95}, can be stated\nin our notation as\n\\begin{eqnarray}\nN_{_{\\rm G}}(x_0,x,y_0,y) \\Big|_{q=2} = {e^{-\\lambda (y-y_0)} \\over 2 x_0\n\\sqrt{\\pi (y-y_0)}} \\left({x \\over x_0}\\right)^{(a+1)\/2} \\, \\exp\n\\left[{-(\\ln x - \\ln x_0)^2 \\over 4 \\, (y-y_0)}\\right]\n\\ ,\n\\label{eq48}\n\\end{eqnarray}\nwhere\n\\begin{equation}\n\\lambda \\equiv {(a+1)^2 \\over 4} + 2 + a + \\theta\n\\ .\n\\label{eq49}\n\\end{equation}\nWe note that the general solution for $N_{_{\\rm G}}$ given by\nequation~(\\ref{eq45}) agrees with equation~(\\ref{eq48}) in the\nlimit $q \\to 2$ as required. Furthermore, the equation~(\\ref{eq45})\nallows $q$ to take on {\\it negative} values if desired, and it is also\napplicable over a broad range of both positive and negative values for\n$a$. Recall that when $a=0$, there are no systematic acceleration or\nloss processes included in the model. Positive (negative) values for $a$\nimply additional systematic acceleration (losses).\n\n\\subsection{Transition to the Stationary Solution}\n\nThe analytical solutions we have obtained for the Green's function\nprovide a complete description of the response of the system to the\nimpulsive injection of monoenergetic particles at any desired time. The\ngenerality of these expressions allows one to obtain the particular\nsolution for the distribution function $f$ associated with any\narbitrary momentum-time source function $S$ using the convolution\ngiven by equation~(\\ref{eq6}). One case of special interest is the\nspectrum resulting from the {\\it continual injection} of monoenergetic\nparticles beginning at time $t=0$, described by the source term\n\\begin{equation}\nS(p,t) = \\cases{\n\\dot N_0 \\, \\delta(p-p_0) \\ , & $t \\ge 0$ \\ , \\cr\n0 \\ , & $t < 0$ \\ , \\cr\n}\n\\label{eq50}\n\\end{equation}\nwhere $\\dot N_0$ denotes the rate of injection of particles with\nmomentum $p_0$ per unit volume per unit time. We assume that no\nparticles are present for $ t < 0$. Combining equations~(\\ref{eq6})\nand (\\ref{eq50}), we find that the time-dependent distribution function\nresulting from monoenergetic particle injection is given by\n\\begin{equation}\nf(p,t) = \\dot N_0 \\int_0^t f_{_{\\rm G}}(p_0,p,t_0,t) \\, dt_0 \\ .\n\\label{eq51}\n\\end{equation}\nBy transforming to the dimensionless variables $x$ and $y$ and employing\nequation~(\\ref{eq16}), we conclude that the particular solution for the\nnumber distribution associated with continual monoenergetic particle injection\ncan be written as\n\\begin{equation}\nN(x,y) \\equiv 4 \\pi \\, m^3 c^3 \\, x^2 \\, f(x,y)\n= {\\dot N_0 \\over D_*} \\int_0^y N_{_{\\rm G}}(x_0,x,y_0,y) \\, dy_0 \\ ,\n\\label{eq52}\n\\end{equation}\nwhere we have used equation~(\\ref{eq13}) to make the substitution $dt_0\n= dy_0 \/ D_*$. Since $N_{_{\\rm G}}$ depends on the time parameters $y$ and\n$y_0$ only through the combination $y-y_0$ (see eqs.~[\\ref{eq45}] and\n[\\ref{eq48}]), it follows that\n\\begin{equation}\nN(x,y) = {\\dot N_0 \\over D_*} \\int_0^y N_{_{\\rm G}}(x_0,x,0,y') \\, dy' \\ .\n\\label{eq53}\n\\end{equation}\nThis is a more convenient form for $N(x,y)$ because $y$ now appears only\nin the upper integration bound.\n\nFor general, finite values of $y$, the time-dependent particular\nsolution for $N(x,y)$ given by equation~(\\ref{eq53}) must be computed\nnumerically by substituting for $N_{_{\\rm G}}$ using either\nequation~(\\ref{eq45}) or (\\ref{eq48}), depending on the value of $q$.\nHowever, as $y \\to \\infty$, the solution rapidly approaches a stationary\nresult representing a balance between injection, acceleration, and\nparticle escape. The form of the stationary solution, called the {\\it\nsteady-state Green's function}, $N^{\\rm G}_{\\rm ss}$, can be obtained by\ndirectly solving the transport equation~(\\ref{eq1}) with $\\partial\nf\/\\partial t=0$ and $S(p,t) = \\dot N_0 \\, \\delta(p-p_0)$. Alternatively,\nthe steady-state Green's function can also be computed by taking the\nlimit of the time-dependent solution, which yields\n\\begin{equation}\nN^{\\rm G}_{\\rm ss}(x) \\equiv \\lim_{y \\to \\infty}\nN(x,y) = {\\dot N_0 \\over D_*} \\int_0^\\infty N_{_{\\rm G}}(x_0,x,0,y')\n\\, dy' \\ .\n\\label{eq54}\n\\end{equation}\nIn the $q < 2$ case, we can substitute for $N_{_{\\rm G}}$ using equation~(\\ref{eq45})\nand evaluate the resulting integral using equation~(6.669.4) from \\citet{gr80}.\nAfter some algebra, we obtain\n\\begin{equation}\nN^{\\rm G}_{\\rm ss}(x) = {\\dot N_0 \\over (2-q) \\, x_0 D_*}\n\\left(x \\over x_0\\right)^{(a+1)\/2} (x x_0)^{(2-q)\/2}\n\\, I_{\\beta-1 \\over 2} \\left(\\sqrt{\\theta} \\, x_{\\rm min}^{2-q} \\over 2-q\\right)\n\\, K_{\\beta-1 \\over 2} \\left(\\sqrt{\\theta} \\, x_{\\rm max}^{2-q} \\over 2-q\\right)\n\\ ,\n\\label{eq55}\n\\end{equation}\nwhere $\\beta = (a+3)\/(2-q)$, and\n\\begin{equation}\nx_{\\rm min} \\equiv \\min(x,x_0) \\ , \\ \\ \\ \\ \\ \\ \\\nx_{\\rm max} \\equiv \\max(x,x_0)\n\\ .\n\\label{eq56}\n\\end{equation}\nLikewise, for the case with $q=2$, we can substitute for $N_{_{\\rm G}}$ using\nequation~(\\ref{eq48}) and then utilize equation~(3.471.9) from \\citet{gr80}\nto conclude that the steady-state solution is given by\n\\begin{equation}\nN^{\\rm G}_{\\rm ss}(x) \\Big|_{q=2} = {\\dot N_0 \\over 2 \\, x_0 D_* \\sqrt{\\lambda}}\n\\left(x \\over x_0\\right)^{(a+1)\/2} \\left(x_{\\rm max} \\over x_{\\rm min}\\right)^{-\\sqrt{\\lambda}}\n\\ ,\n\\label{eq57}\n\\end{equation}\nwhere $\\lambda$ is defined by equation~(\\ref{eq49}). The steady-state\nsolutions given by equations~(\\ref{eq55}) and (\\ref{eq57}) agree with\nthe results obtained by directly solving the transport equation. Due to\nthe asymptotic behavior of the Bessel $K_\\nu(z)$ function for large $z$,\nequation~(\\ref{eq55}) indicates that $N^{\\rm G}_{\\rm ss}$ exhibits an\nexponential cutoff at high energies when $q < 2$ \\citep{as70}. This\ncorresponds physically to the fact that the escape timescale decreases\nwith increasing particle momentum in this case. However, when $q=2$\n(eq.~[\\ref{eq57}]), the spectrum displays a pure power-law behavior at\nhigh energies because the escape timescale is independent of the\nparticle momentum. Specific examples illustrating these behaviors will be\npresented in \\S~4.\n\n\\subsection{Escaping Particle Distribution}\n\nThe various expressions we have obtained for the Green's function\n$N_{_{\\rm G}}$ describe the momentum distribution of the particles remaining\nin the plasma at time $t$ after injection occurring at time $t_0$.\nHowever, since our model incorporates a physically realistic,\nmomentum-dependent escape timescale $t_{\\rm esc}(p)$ given by\nequation~(\\ref{eq7}), it is also quite interesting to compute the\nspectrum of the {\\it escaping} particles, which may form an energetic\noutflow capable of producing observable radiation. In general, the\nnumber distribution of the escaping particles, $\\dot N^{\\rm esc}(x,y)$,\nis related to the current distribution of particles in the plasma, $N(x,y)$,\nvia\n\\begin{equation}\n\\dot N^{\\rm esc}(x,y) \\equiv t_{\\rm esc}^{-1} \\, N(x,y)\n= \\theta D_* \\ x^{2-q} \\, N(x,y)\n\\ ,\n\\label{eq58}\n\\end{equation}\nwhere we have used equations~(\\ref{eq7}) and (\\ref{eq14}) to obtain the\nfinal result. The quantity $\\dot N^{\\rm esc}(x,y)\\,dx$ represents the number\nof particles escaping per unit volume per unit time with dimensionless\nmomenta between $x$ and $x+dx$.\n\nAn important special case is the evolution of the escaping particle\nspectrum resulting from {\\it impulsive monoenergetic injection} at\ndimensionless time $y=y_0$. In this application, equation~(\\ref{eq58})\ngives the Green's function number spectrum for the escaping particles,\ndefined by\n\\begin{equation}\n\\dotN_{_{\\rm G}}^{\\rm esc}(x_0,x,y_0,y)\n\\equiv \\theta D_* \\ x^{2-q} \\, N_{_{\\rm G}}(x_0,x,y_0,y)\n\\ ,\n\\label{eq59}\n\\end{equation}\nwhere $N_{_{\\rm G}}$ is evaluated using either equation~(\\ref{eq45}) or\n(\\ref{eq48}) depending on the value of $q$. We can use this expression\nfor $\\dotN_{_{\\rm G}}^{\\rm esc}$ to compute the {\\it total} escaping number\ndistribution resulting from an impulsive flare occurring at $t=0$ by\nintegrating the escaping distribution with respect to time, which yields\n\\begin{equation}\n\\DeltaN_{_{\\rm G}}^{\\rm esc}(x) \\equiv \\int_0^\\infty \\dotN_{_{\\rm G}}^{\\rm esc} \\, dt\n= \\theta \\, x^{2-q} \\int_0^\\infty N_{_{\\rm G}}(x_0,x,0,y') \\, dy'\n\\ ,\n\\label{eq60}\n\\end{equation}\nwhere we have used equation~(\\ref{eq59}) and taken advantage of the fact\nthat $N_{_{\\rm G}}$ depends on $y$ and $y_0$ only through the difference $y-y_0$\n(cf. eq.~[\\ref{eq53}]).\nBy comparing equations~(\\ref{eq54}) and\n(\\ref{eq60}), we deduce that\n\\begin{equation}\n\\DeltaN_{_{\\rm G}}^{\\rm esc}(x) = {\\theta D_* \\over \\dot N_0}\n\\ x^{2-q} \\, N^{\\rm G}_{\\rm ss}(x)\n\\ ,\n\\label{eq61}\n\\end{equation}\nso that the total escaping spectrum is simply proportional to the steady-state\nGreen's function resulting from continual injection, as expected. The process\nconsidered here corresponds to the injection of a single particle at time $t=0$,\nand therefore we find that the normalization of $\\DeltaN_{_{\\rm G}}^{\\rm esc}$ is\ngiven by\n\\begin{equation}\n\\int_0^\\infty \\DeltaN_{_{\\rm G}}^{\\rm esc}(x) \\, dx = 1\n\\ ,\n\\label{eq62}\n\\end{equation}\nwhich provides a useful check on the numerical results.\n\nAnother interesting example is the case of continual monoenergetic\nparticle injection commencing at time $t=0$, described by the source\nterm given by equation~(\\ref{eq50}). The time-dependent buildup of the\nescaping particle spectrum $\\dot N^{\\rm esc}(x,y)$ in this scenario can\nbe analyzed by using equation~(\\ref{eq53}) to substitute for $N(x,y)$ in\nequation~(\\ref{eq58}), which yields\n\\begin{equation}\n\\dot N^{\\rm esc}(x,y)\n= \\dot N_0 \\, \\theta \\ x^{2-q} \\int_0^y N_{_{\\rm G}}(x_0,x,0,y') \\, dy'\n\\ .\n\\label{eq63}\n\\end{equation}\nIn the limit $y \\to \\infty$, the escaping particle spectrum approaches the\nsteady-state result\n\\begin{equation}\n\\dot N_{\\rm ss}^{\\rm esc}(x) \\equiv \\lim_{y \\to \\infty} \\dot N^{\\rm esc}(x,y)\n= \\dot N_0 \\, \\theta \\ x^{2-q} \\int_0^\\infty N_{_{\\rm G}}(x_0,x,0,y') \\, dy'\n\\ .\n\\label{eq64}\n\\end{equation}\nBy comparing equations~(\\ref{eq54}), (\\ref{eq61}), and (\\ref{eq64}), we\nconclude that\n\\begin{equation}\n\\dot N_{\\rm ss}^{\\rm esc}(x)\n= \\theta D_* \\ x^{2-q} \\, N^{\\rm G}_{\\rm ss}(x)\n= \\dot N_0 \\ \\DeltaN_{_{\\rm G}}^{\\rm esc}(x)\n\\ .\n\\label{eq65}\n\\end{equation}\nThe final result can be combined with equation~(\\ref{eq62}) to show that\nthe escaping spectrum satisfies the normalization relation\n\\begin{equation}\n\\int_0^\\infty \\dot N_{\\rm ss}^{\\rm esc}(x) \\, dx = \\dot N_0 \\ ,\n\\label{eq66}\n\\end{equation}\nas expected for the case of continual steady-state injection.\n\nNote that analytical expressions for the steady-state escaping spectrum $\\dot\nN_{\\rm ss}^{\\rm esc}(x)$ can be obtained by substituting for $N^{\\rm G}_{\\rm\nss}(x)$ in equation~(\\ref{eq65}) using either equation~(\\ref{eq55}) or\n(\\ref{eq57}), depending on the value of $q$. We obtain\n\\begin{equation}\n\\dot N^{\\rm esc}_{\\rm ss}(x) = {\\dot N_0 \\, \\theta \\, x^{2-q} \\over (2-q) \\, x_0}\n\\left(x \\over x_0\\right)^{(a+1)\/2} (x x_0)^{(2-q)\/2}\n\\, I_{\\beta-1 \\over 2} \\left(\\sqrt{\\theta} \\, x_{\\rm min}^{2-q} \\over 2-q\\right)\n\\, K_{\\beta-1 \\over 2} \\left(\\sqrt{\\theta} \\, x_{\\rm max}^{2-q} \\over 2-q\\right)\n\\ ,\n\\label{eq67}\n\\end{equation}\nfor $q < 2$, and\n\\begin{equation}\n\\dot N^{\\rm esc}_{\\rm ss}(x) = {\\dot N_0 \\, \\theta \\, x^{2-q} \\over\n2 \\, x_0 \\sqrt{\\lambda}}\n\\left(x \\over x_0\\right)^{(a+1)\/2} \\left(x_{\\rm max} \\over x_{\\rm min}\\right)^{-\\sqrt{\\lambda}}\n\\ ,\n\\label{eq68}\n\\end{equation}\nfor $q = 2$.\n\n\\section{RESULTS}\n\nThe new result we have derived for the secular Green's function\n(eq.~[\\ref{eq45}]) displays a rich behavior through its complex\ndependence on momentum, time, and the dimensionless parameters $q$, $a$,\nand $\\theta$, which are related to the fundamental physical transport\ncoefficients $D_*$, $A_*$, and $t_*$ via equations~(\\ref{eq14}). Here we\npresent several example calculations in order to illustrate the utility\nof the new solution. Detailed applications to astrophysical situations,\nincluding active galaxies and $\\gamma$-ray bursts, will be presented in\nsubsequent papers.\n\n\\begin{figure}[t]\n\\includegraphics[width=6in]{f2.eps}\n\\vskip-0.6in\n\\caption{Green's function solutions to the time-dependent stochastic\nparticle acceleration equation for $x_0=1$. The left panels depict the\nimpulsive-injection solution, $N_{_{\\rm G}}$ (eqs.~[\\ref{eq45}] and\n[\\ref{eq48}]), and the right panels illustrate the response to uniform,\ncontinuous injection, $N$ (eq.~[\\ref{eq53}]), with $\\dot\nN_0=D_*=1$. Note that $N$ approaches the corresponding steady-state\nsolution (eqs.~[\\ref{eq55}] and [\\ref{eq57}]) as $y$ increases. The\nindices of the wave turbulence spectra are indicated, with $q = 2$ for\nhard-sphere scattering, $q = 5\/3$ for a Kolmogorov cascade, and $q =\n3\/2$ for a Kraichnan cascade.}\n\\label{fig2}\n\\end{figure}\n\nThe panels on the left-hand side of Figure~2 depict the time-dependent\nGreen's function solution, $N_{_{\\rm G}}$, describing the evolution of the\nparticle distribution in the plasma resulting from {\\it impulsive}\nmonoenergetic injection at $y=0$. Results are presented for the\nhard-sphere scattering case ($q = 2$), computed using\nequation~(\\ref{eq48}), and for the Kolmogorov ($q = 5\/3$) and Kraichnan\n($q = 3\/2$) cases, evaluated using equation~(\\ref{eq45}). The only\nacceleration mechanism considered here is the stochastic acceleration\nassociated with the second-order Fermi process, and therefore we set $a\n= 0$. The escape parameter $\\theta$ is set equal to unity, so that the\ntimescale for escape is comparable to the diffusion timescale. As the\nwave turbulence spectrum becomes steeper (i.e., as $q$ increases), a\nlarger fraction of the turbulence energy is contained in waves with long\nwavelengths, which interact resonantly with higher energy particles.\nSteeper turbulence spectra therefore produce harder particle\ndistributions, as can be confirmed in the plots. Consequently we\nconclude that an ensemble of hard sphere scattering centers is more\neffective at accelerating nonthermal particles compared with a\nKolmogorov wave spectrum, which in turn is more effective than a\nKraichnan spectrum.\n\nThe panels on the right-hand side of Figure~2 illustrate the buildup of\nthe particle spectrum in the plasma, $N(x,y)$, due to {\\it continual}\nmonoenergetic injection beginning at $y=0$, computed by evaluating\nnumerically the integral in equation~(\\ref{eq53}). We have set $a=0$ and\ntherefore the acceleration is purely stochastic. As $y \\to \\infty$, the\nspectrum approaches the steady-state form given by equation~(\\ref{eq55})\nfor $q<2$ or by equation~(\\ref{eq57}) for $q=2$. In the hard-sphere case\n($q=2$), the particle spectrum displays a power-law shape at high\nenergies in agreement with equation~(\\ref{eq57}). However, when $q < 2$,\nparticle escape dominates over acceleration at high energies, and\ntherefore the steady-state distribution is truncated, even in the\nabsence of radiative losses. This effect produces the quasi-exponential\nturnovers exhibited by the stationary spectra when $q = 5\/3$ and $q =\n3\/2$. Particle escape in these cases mimics energy losses due to, for\nexample, synchrotron emission from electrons.\n\n\\begin{figure}[t]\n\\hskip0.2in\\plottwo{f3a.eps}{f3b.eps}\n\\hskip0.5truein\n\\caption{Evolution of the particle distribution in the plasma resulting\nfrom monoenergetic injection with $q = 5\/3$, $a = 0$, $\\theta = 0.1$,\nand $x_0=1$. Panel~({\\it a}) depicts the Green's function, $N_{_{\\rm G}}$,\nresulting from impulsive monoenergetic injection (eq.~[\\ref{eq45}]), and\npanel ({\\it b}) illustrates the response to continual monoenergetic\ninjection (eq.~[\\ref{eq53}]) and the corresponding steady-state solution\n(eq.~[\\ref{eq55}]) for $\\dot N_0=D_*=1$.}\n\\label{fig3}\n\\end{figure}\n\nFigures~3 and 4 illustrate the effects of varying the values of the\nescape parameter $\\theta$ and the systematic acceleration\/loss parameter\n$a$ for the $q = 5\/3$ case. In Figure~3, we set $a=0$ and $\\theta =\n0.1$, which represents a particle escape timescale that is an order of\nmagnitude larger than the corresponding case depicted in Figure~2. The\nlonger escape time allows the particles to be accelerated to higher mean\nenergies before diffusing out of the plasma, and the decay of the\nparticle density at late times takes place much more slowly than for\nlarger values of $\\theta$. The hardening of the spectra causes the\nquasi-exponential cutoffs to move to higher energies, and the same\neffect can also be noted in the corresponding stationary solutions. The\ncalculations represented in Figure~4 include additional systematic\nacceleration processes that are modeled by setting $a = 2$. The enhanced\nparticle acceleration further hardens the spectrum and shifts the cutoff\nto even higher energies. Although the slope of the low-energy particle\ndistribution ($x < x_0$) is not altered much when only $\\theta$ is\nvaried (see Figs.~2 and 3), this slope becomes significantly steeper\nwhen $a$ is increased, as indicated in Figure~4.\n\nIn the left-hand panels of Figures~5 and 6 we plot the time-dependent\nsolution for the Green's function describing the escaping particles,\n$\\dotN_{_{\\rm G}}^{\\rm esc}$, resulting from {\\it impulsive} monoenergetic\nparticle injection (eq.~[\\ref{eq59}]) when $q = 5\/3$. The right-hand\npanels illustrate the buildup of the escaping spectrum, $\\dot N^{\\rm\nesc}$ (eq.~[\\ref{eq63}]), resulting from {\\it continual} injection\nbeginning at $y=0$. Note the transition to the steady-state solution,\n$\\dot N_{\\rm ss}^{\\rm esc}$ (eq.~[\\ref{eq64}]), as $y \\to \\infty$. In\nFigure~5 we set $\\theta=0.1$ and $a=0$, and in Figure~6 we set\n$\\theta=0.1$ and $a=2$. Included for comparison are the corresponding\nspectra describing the particle distributions in the plasma at the same\nvalues of $y$. We point out that the escaping particle spectra are\nsignificantly harder than the in situ distributions, which reflects the\npreferential escape of the high-energy particles resulting from the\nmomentum dependence of the escape timescale when $q<2$ (see\neq.~[\\ref{eq7}]).\n\n\\begin{figure}[t]\n\\hskip0.2in\\plottwo{f4a.eps}{f4b.eps}\n\\caption{Same as Fig.~3, except $a = 2$.}\n\\label{fig4}\n\\end{figure}\n\n\\section{DISCUSSION AND SUMMARY}\n\nWe have derived new, closed-form solutions (eqs.~[\\ref{eq45}] and\n[{\\ref{eq67}]}) for the time-dependent Green's function representing the\nstochastic acceleration of relativistic ions interacting with MHD waves.\nThe analytical results we have obtained describe the time-dependent\ndistributions for both the accelerated (in situ) and the escaping\nparticles. The Fokker-Planck transport equation considered here includes\nmomentum diffusion with coefficient $D(p) \\propto p^q$, particle escape\nwith mean timescale $t_{\\rm esc} \\propto p^{q-2}$, and additional\nsystematic acceleration\/losses with a rate proportional to $p^{q-1}$,\nwhere $p$ is the particle momentum and $q$ is the index of the wave\nturbulence spectrum. This specific scenario describes the resonant\ninteraction of particles with fast-mode and Alfv\\'en waves, which is one\nof the fundamental acceleration processes in high-energy astrophysics\n\\citep[e.g.,][]{dml96}.\n\nThe new analytical result complements the work of \\citet{pp95} since it\nis applicable for any $q < 2$ provided $s_{\\rm pp} = q-2$, where $s_{\\rm\npp}$ is the power-law index used by these authors to describe the\nmomentum dependence of the escape timescale. The most closely related\nprevious result is given by their equation~(59), which treats the case\nwith $q \\ne 2$ and $s_{\\rm pp}=0$, corresponding to a\nmomentum-independent escape timescale. Our analytical solution\n(eq.~[\\ref{eq45}]) agrees with theirs in the limit $q \\to 2$, as\nexpected. The general features of our new solution were discussed in\n\\S~4, where it was demonstrated that increasingly hard particle spectra\nresult from larger values of the wave index $q$, smaller values of the\nescape parameter $\\theta$, and larger values of the systematic\nacceleration parameter $a$. The rich behavior of the solution as a\nfunction of momentum, time, and the parameters $q$, $\\theta$, and $a$\nprovides useful physical insight into the nature of the coupled\nenergetic\/spatial particle transport in astrophysical plasmas.\n\n\\begin{figure}[t]\n\\hskip0.2in\\plottwo{f5a.eps}{f5b.eps}\n\\caption{Evolution of the particle distribution resulting from\nmonoenergetic injection with $q = 5\/3$, $a = 0$, $\\theta = 0.1$, $x_0 =\n1$, $\\dot N_0=1$, and $D_*=1$. Panel~({\\it a}) treats the case of\nimpulsive injection, with the thin lines representing the particle\ndistribution in the plasma (eq.~[\\ref{eq45}]) and the heavy lines denoting\nthe escaping particle spectrum (eq.~[\\ref{eq59}]) in units of $\\theta$.\nPanel~({\\it b}) illustrates the response to continual monoenergetic\ninjection for the particle distribution in the plasma (eqs.~[\\ref{eq53}] and\n[\\ref{eq55}]; thin lines) and for the escaping particle spectrum\n(eqs.~[\\ref{eq63}] and [\\ref{eq65}]; heavy lines).\nSee the discussion in the text.}\n\\label{fig5}\n\\end{figure}\n\n\\begin{figure}[t]\n\\hskip0.2in\\plottwo{f6a.eps}{f6b.eps}\n\\caption{Same as Fig.~5, except $a = 2$.}\n\\label{fig6}\n\\end{figure}\n\nThe solutions presented here can be used to describe the acceleration\nand transport of relativistic ions in astrophysical environments in\nwhich the turbulence spectrum is very poorly known and can be\napproximated by a power law, such as $\\gamma$-ray bursts, active\ngalaxies, magnetized coronae around black holes, and the intergalactic\nmedium in clusters of galaxies. For example, the hard X-ray emission\nfrom black-hole jet sources such as Cygnus X-1 \\citep{mal06} and the\nmicroquasar LS 5039 \\citep{aha05} could be powered by the stochastic\nacceleration of ions in a black-hole accretion disk corona that\nsubsequently escape and interact with surrounding material. In this\nscenario, persistent acceleration of monoenergetic particles injected\ninto the corona, or flaring episodes averaged over a sufficiently long\ntime, would produce a time-averaged escaping distribution of\nrelativistic protons given by equation~(\\ref{eq67}) for $q < 2$.\nAssuming only stochastic acceleration, so that $a = 0$, the number\ndistribution of escaping particles with $x \\geq x_0$ takes the form\n\\begin{equation}\n\\dot N_{\\rm ss}^{\\rm esc}(x) \n\\propto x^{(7-3q)\/2} K_{{\\beta - 1\\over 2}}\\left( {\\sqrt{\\theta} \nx^{2-q} \\over 2-q}\\right) \\propto \\cases{\nx^{3-2q} \\ , & $x \\ll \\left( {2-q\\over \\sqrt{\\theta}}\\right)^{1\/(2-q)}$ \\ \\cr\nx^{(5-2q)\/2}\\exp\\left(-{\\sqrt{\\theta}x^{2-q}\\over 2-q}\\right) \\ , & \n$x \\gg \\left({2-q\\over \\sqrt{\\theta}}\\right)^{1\/(2-q)}$ \\ \\cr\n}.\n\\label{eq69}\n\\end{equation}\nWhen the escaping hadrons collide with ambient gas or stellar wind\nmaterial, they would generate X-rays and $\\gamma$-rays via a pion\nproduction cascade with a very hard spectrum leading up to a\nquasi-exponential cutoff. The {\\it Gamma ray Large Area Space Telescope}\n(GLAST) telescope will be able to provide detailed spectra from Galactic\nblack-hole sources and unidentified EGRET $\\gamma$-ray sources to test\nfor the existence of this hard component. Our new analytical solution\ncan also be used to model the variability of radiation from ions\naccelerated in the accretion-disk coronae of Seyfert galaxies by\nchanging the level of turbulence. Flaring $\\sim 100\\,$MeV -- GeV\nradiation could be weakly detected by GLAST as a consequence of this\nprocess. Additional applications of our work include studies of the\nstochastic acceleration of relativistic cosmic rays in $\\gamma$-ray\nbursts \\citep{dh01}.\n\nThe results we have obtained describe both the momentum distribution of\nthe particles in the plasma, and the momentum distribution of the\nparticles that escape to form energetic outflows. Since the solutions do\nnot include inverse-Compton or synchrotron losses, which are usually\nimportant for energetic electrons, they are primarily applicable to\ncases involving ion acceleration. In order to treat the acceleration of\nrelativistic electrons, a generalized calculation including both\nstochastic particle acceleration and radiative losses is desirable, and\nwe are currently working to extend the analytical model presented here\nto incorporate these effects. Beyond the direct utility of the new\nanalytical solutions for probing the nature of particle acceleration in\nastrophysical sources, we note that they are also useful for\nbenchmarking numerical simulations.\n\n\\acknowledgements\nT.~L. is funded through NASA {\\it GLAST} Science Investigation No.~DPR-S-1563-Y.\nThe work of C.~D.~D. and P.~A.~B. is supported by the Office of Naval Research.\nThe authors also acknowledge the useful comments provided by the anonymous\nreferee.\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzqkol b/data_all_eng_slimpj/shuffled/split2/finalzzqkol new file mode 100644 index 0000000000000000000000000000000000000000..a0f0e4671a086068cf17094b1c2d62a47825bf40 --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzqkol @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\n Phase retrieval aims to reconstruct a complex-valued signal from its intensity-only measurements. It is a crucial problem in crystallography, optical imaging, astronomy, X-ray, electronic imaging, etc., because most existing measurement systems can only detect the relative low-frequency magnitude information of the wave fields. The study of the phase retrieval problems originated in the $1970$s. The researchers of the optics community developed many reconstruction algorithms \\cite{Gerchberg1972APA, Fienup:82, rodenburg2008ptychography}. \n \n Mathematically, the phase retrieval problem can be described by the following equation:\n \\begin{equation}\n \t\\label{CDP}\n \t\\begin{aligned}\n \t\t\\text { Find } \\mathbf{x} \\in \\mathbb{C}^{N} \\quad \\text { s.t. } \\ \\mathbf{y}_{i}=\\left|\\mathcal{F}\\left(\\mathbf{I}_{i} \\circ \\mathbf{x}\\right)\\right|^{2}, i=1, \\ldots, M,\n \t\\end{aligned}\n \\end{equation}\n where $\\mathbf{x}\\in \\mathbb{C}^N$ is the complex-valued signal of interest, and $ \\mathbf{y}$ is its Fourier intensity measurements; $\\circ $ and $\\mathcal{F}$ refer to elementwise multiplication and Fourier transform operator, respectively. The pre-defined optical masks $ \\mathbf{I} $ provide the constraints to reduce the ill-posedness of the problem. It is implemented in many different ways. For instance, early-stage phase retrieval algorithms treat the support of the signal as the optical mask \\cite{Fienup:82}. Specifically, $|\\mathbf{x}_i| \\neq 0$ when $i \\in \\mathcal{S}$, where $\\mathcal{S}$ denotes the support of the signal. However, these early-stage algorithms cannot guarantee the convergence of the optimization process nor provide globally optimal solutions. Recently, pre-defined random optical masks have been used as the constraints to improve the reconstruction performance \\cite{Anand07, Horisaki2014SingleshotPI, candes2015phase, YE2022106808}. The random masks can be implemented using a spatial light modulator (SLM) or digital micromirror device (DMD) \\cite{Horisaki2014SingleshotPI, Zheng2017DigitalMD}. Although the usage of the random masks can lead to better reconstruction performance, there are several disadvantages. The cost of DMD and SLM is one concern, the error due to the global gray-scale phase mismatch and spatial non-uniformity from the devices is another. Besides, it was empirically shown that about $ 4 - 6 $ measurements are required for exact recovery with random binary masks \\cite{candes2015code, candes2015phase}. It increases the measurement time and can affect the quality of the reconstructed image, particularly for dynamic objects.\n \n Recently, data-driven methods, like deep neural networks, have been widely applied to explore the data distributions through extensive training samples \\cite{LeCun_2015}. Among them, the so-called Convolutional Neural Network (CNN) that exploits the features based on convolution kernels has been successfully adopted in image processing tasks, like image restoration \\cite{He_2016_CVPR, zhang2017beyond, wang2020lightening}. Besides, an important technique named attention mechanism that mimics cognitive attention became popular in the CNN structure to enhance some parts of the feature maps and thus improve the performance \\cite{attention_17}. Since CNN can learn discriminative representations from large training datasets, researchers have spent effort in using CNN for solving phase retrieval problems \\cite{deng_li_goy_kang_barbastathis_2020, White:19, Shi:19, Sinha17}. For instance, Chen et al. reconstructed signals from Fresnel diffraction patterns with CNN and the contrast transfer function \\cite{Bai:19}. Kumar et al. proposed the deep unrolling networks that solved the iterative phase retrieval problem with CNN \\cite{Kumar:21}. Besides, the conditional generative adversarial network (GAN), a kind of data-driven generative method based on specific conditions \\cite{cond_gan_14}, was adopted for reconstructing the phase information from Fourier intensity measurements \\cite{uelwer2021phase}. Wu et al. also proposed a CNN to reconstruct the crystal structures with coherent X-ray diffraction patterns \\cite{Wu_cw5029, Wu_2021}. The above methods adopt supervised learning that trains the networks with the paired ground-truth data. In recent years, training with unpaired or unsupervised (no ground-truth data) datasets has been a popular topic. For example, Fus et al. proposed the unsupervised learning for in-line holography phase retrieval \\cite{Fus_18}, while Zhang et al. developed the GAN-based phase retrieval with unpaired datasets \\cite{Zhang_21}. While the above approaches have achieved some success, phase retrieval with Fraunhofer diffraction patterns (i.e., Fourier intensity) for general applications is still a challenge. It is partly because of the large domain discrepancy between the Fourier plane and image plane; the large variation in data distributions of general applications also introduces much difficulty in designing a model structure that is optimal for different applications. In this paper, we propose a single-shot maskless phase retrieval method named \\textit{SiPRNet} that is benefited from the advanced deep learning technology. Compared with the traditional algorithms that need multiple intensity measurements, the proposed method only collects one Fourier intensity measurement for reconstruction. Besides, no optical mask is needed in the proposed method; thus, the cost for expensive SLM\/DMD devices is saved. The proposed \\textit{SiPRNet} is designed to work with general image datasets without assuming a specific application domain. It has demonstrated its generality by achieving state-of-the-art performance on three different datasets, as shown in our experimental results. It significantly outperforms the existing optimization-based and learning-based phase retrieval methods.\n \n The contributions of our method are: (1) We propose a novel CNN structure - \\textit{SiPRNet} that utilizes advanced deep learning techniques for single-shot maskless phase retrieval. It can reliably reconstruct the original signal based on only one Fourier intensity measurement without using optical masks. (2) We propose a Multi-Layer Perception (MLP) block to extract global representations from the intensity measurement. In addition, we propose to embed a dropout layer into the MLP block to constrain the behavior of the fully connected layers, which promotes the MLP to extract global representation. (3) We propose to use residual learning and self-attention mechanism in the Up-sampling and Reconstruction blocks, which improves the quality of the reconstructed signals. (4) We verify the proposed method on an optical platform to demonstrate its practicality. The source codes of the proposed method have been released on Github: \\url{https:\/\/github.com\/Qiustander\/SiPRNet}.\n \n The rest of this paper is organized as follows. Section \\ref{Sec:Method} introduces the defocused phase retrieval system we developed in this study and the proposed end-to-end deep-learning phase retrieval network \\textit{SiPRNet}. Section \\ref{Sec:simulation} presents the model analysis and simulation results. Section \\ref{Sec:experiment} provides the experimental results on an optical system for validating the performance of the proposed method. \n \n\\section{Proposed Method} \n\\label{Sec:Method}\n\n\\subsection{Phase Retrieval System} \\label{Sec:PRsystem}\n\nAs indicated in Eq. \\eqref{CDP}, a phase retrieval system aims to reconstruct a complex-valued signal $ \\mathbf{x} $ from its Fourier intensity measurements. The proposed method does not require multiple measurements or optical masks, so $ M=1 $ and all the values of the mask $ \\mathbf{I} $ are equal to one in Eq. \\eqref{CDP}. In this case, the original complex-valued signal $ \\mathbf{x} $ is reconstructed only from its Fourier intensity. This seemingly impossible task is known to have a solution, although with ambiguities. From \\cite{hayes1982reconstruction}, it is known that $ \\mathbf{x} $ can be uniquely defined, except for trivial ambiguities, by its Fourier intensity, with an oversampling factor over $ 2 $ in each dimension, if $ \\mathbf{x} $ has a finite support and is non-symmetric. The above has an important implication to the optical system required for the proposed phase retrieval method. Fig. \\ref{fig:opticalpath} shows the system we have constructed to implement the proposed method. In the system, we increase the resolution of the CMOS camera to make sure that the captured Fourier intensity image has a sampling rate at least two times in each dimension higher than that of the object image $ \\mathbf{x} $. As a result, the number of samples in the zero diffraction order of the Fourier intensity image is $ 762\\times762 $, while the number of samples of $ \\mathbf{x} $ is $ 128\\times128 $. We then extract the central $ 128\\times128 $ pixels of the Fourier intensity measurements for feeding to the proposed CNN model for phase retrieval. \n\n\\begin{figure}[htb] \n\t\\centering\\includegraphics[width=0.8\\linewidth]{Figures\/Fig_Optical_Path.pdf}\n\t\\caption{Optical path of the defocus phase retrieval system. $\\mathbf{x}$ represents the complex-valued object, and $\\mathbf{y}$ is its Fourier intensity measurement; $\\mathbf{X} = \\mathcal{F}\\left(\\mathbf{I}_{i} \\circ \\mathbf{x}\\right)$, $\\mathbf{H}$ and $\\ast$ denote the Fourier plane in complex-valued form, the defocus function generated via Fresnel diffraction and the convolution operation, respectively. }\n\t\\label{fig:opticalpath}\n\t\n\\end{figure}\n\nWhen capturing Fourier intensity images, it is common to have the saturation problem. It is because most structured signals have the energy concentrated in the low-frequency areas, in particular the zero frequency. Due to the limited analog-to-digital resolution of standard scientific cameras ($14$ - $16$ bits), the central regions of the Fourier intensity measurement are often seriously saturated. To lessen the effect of the saturation problem, we convolve the Fourier intensity measurement with a defocus function $ \\mathbf{H} $ to reduce the dynamic range \\cite{goodman2017introduction}. It can be implemented by moving the camera beyond the Fourier plane. Fig. \\ref{fig:opticalpath} shows a typical optical path for implementing the defocus phase retrieval system. Fig. \\ref{fig:compvisual} shows some examples of the Fourier intensity measurements obtained in our experiments. \n\nBesides the lens-based system that we have implemented in this work, we can also find lensless phase retrieval systems \\cite{Sinha17}. Different from the lens-based system that reconstructs images from their Fraunhofer diffraction patterns, the lensless system reconstructs images based on their Fresnel diffraction patterns \\cite{Sinha17}. Compared with the lens-based system, the lensless system requires a much longer distance between the sensors and objects. For instance, in \\cite{Sinha17}, the distance between the sensor and object is reported as $37.5 cm$ or longer. It would lead to severe aberrations and energy loss. By using a lens, the lens-based system can generate high-quality Fraunhofer diffraction patterns at a short propagation distance between the lens and sensor ($3 cm$ or less), thus avoiding the drawbacks stated above.\n\n\n\\subsection{Deep Learning Model: \\textit{SiPRNet}}\n\nWith the oversampling constraint as mentioned above, it is possible to reconstruct a complex-valued signal (with trivial ambiguities) from only its Fourier intensity without other constraints (such as optical masks). However, such a reconstruction problem is highly non-convex that is difficult to solve with traditional optimization methods. Due to their non-linear characteristic, deep neural networks have been successfully applied to many non-convex optimization problems. In this paper, we propose to regard the single-shot maskless phase retrieval challenge as an end-to-end learning task, and formulate the retrieval process directly by a deep neural network model (denoted as $P(\\mathbf{y};\\theta)$) as:\n\n\\begin{dmath}\n\t\\hat{\\mathbf{x}} = P(\\mathbf{y};\\theta),\n\\end{dmath}\nwhere $\\theta$ represents the learnable parameters of the deep learning model. The training target is to minimize the error between the predicted $P(\\mathbf{y};\\theta)$ and the ground-truth complex-valued image $ \\mathbf{x} $ regularized with the total variation (TV) norm $ TV(\\hat{\\mathbf{x}}) $. The training process can be written as:\n\n\\begin{equation}\n\tP(y;\\theta) = \\mathop{argmin}\\limits_{\\hat{\\mathbf{x}}} \\left( \\left\\| \\mathbf{x}-\\hat{\\mathbf{x}} \\right\\|_1 +\\lambda TV\\left(\\hat{\\mathbf{x}}\\right)\\right), \n\\end{equation}\nwhere $\\left\\| \\cdot \\right\\|_1$ denotes the $\\ell_1$-norm distance and $ \\lambda $ is a constant determined empirically.\n\nFig. \\ref{fig:structure} shows the overall structure of the proposed phase retrieval network: \\textit{SiPRNet}. It mainly consists of two parts: feature extraction and phase reconstruction. It takes the Fourier intensity measurement $ \\mathbf{y} $ as input and reconstruct the complex-valued image $ \\mathbf{x} $ as its output.\n\n\\begin{figure} [htb]\n\t\\centering\n\t\\includegraphics[width=\\linewidth]{Figures\/Fig_Network_Structure.pdf}\n\t\\caption{Structure of the proposed network that mainly consists of two parts: Feature Extraction and Phase Reconstruction.}\n\t\\label{fig:structure}\n\\end{figure}\n\n\\subsubsection{\\textbf{Feature Extraction}}\nTraditional learning-based approaches often use a CNN as the front end to extract the features of the Fourier intensity measurement. However, the convolution operation of CNN shares the parameters (e.g., weights and bias) for all data of the measurement but ignores their different characteristics. Therefore, we propose to use a new Multi-Layer Perceptron (MLP) block to fuse the information of different frequency data as shown in Fig. \\ref{fig:structure}(a). The MLP block consists of three Fully Connected (FC) and one Dropout layer. The input Fourier measurement is firstly reshaped from the size of $128\\times 128\\times 1$ to $1\\times1\\times16,384$, and then fed to the first FC layer of the MLP block. Each FC layer has $1,024$ neurons for extracting the features required for reconstructing the original complex-valued images. We propose to embed a Dropout layer after the first FC layer. The Dropout layer randomly deactivates a small number (empirically $20\\%$) of neurons at the training stage. At the inference stage, all neurons are activated that forms an ensemble to reduce the risk of overfitting \\cite{srivastava2014dropout}. Our design with embedding a Dropout layer after the first FC layer constrains the behavior of the FC layers so that the features extracted by the MLP block do not depend on a particular input frequency measurement. The effectiveness of the dropout layer has been evaluated and shown in Section \\ref{Sec:ablationstudy}. Fig. \\ref{fig:structure} shows an example of the original complex-valued image (real and imaginary), the Fourier intensity measurement, and the features extracted by the MLP block. It can be seen that the MLP block has performed the domain conversion. The extracted feature image contains the features of the original complex-valued image. A decoder is needed to reconstruct the real and imaginary images from the extracted features. \n\n\n\\subsubsection{\\textbf{Phase Reconstruction}}\nThe phase reconstruction part contains two Up-sampling and Reconstruction (UR) blocks and one Post-processing block, as shown in Fig. \\ref{fig:structure}(b) and (c). The UR block is responsible for reconstructing the complex-valued image from the extracted features and progressively upsampling it back to the original size. The basic processing unit of the UR block is the convolution block that has three components: convolutional, instance normalization, and Parametric Rectified Linear Unit (PReLU) layers. The convolutional layer contains $3\\times 3$ learnable filters that aim to implicitly approximate the nonlinear backward mapping from the extracted features to the desired complex-valued image $ \\mathbf{x} $. The instance normalization and PReLU layers are non-linear functions that increase the non-linear capacity and learning ability of the deep learning model. \n\nA self-attention unit is followed after each convolutional block. Recently, the self-attention technique has been popularly used in image reconstruction tasks, e.g., the scattering problem \\cite{non-local, Wang:21}. Traditional CNN using small convolution kernels can only extract the local correlation of an image. However, natural images are often structural that exhibit global correlation. Self-attention is an effective way to exploit the global correlation in images. The structure of the self-attention layer used in the UR block is similar to that in \\cite{non-local}, as shown in Fig. \\ref{fig:attn}. At the first stage, the input features are fed to three $1\\times1$ convolutional layers ($\\theta, \\phi$ and $g$). The outputs are denoted as key, query, and value, respectively, all with the same dimensions. Next, the cross-correlation matrix is computed between the key and query. It represents the dependency between any two pixels of the input features. The softmax function is then performed to normalize the correlation scores and build an attention map. Large values in the attention map denote strong correlations between the features concerned. Since the correlation is computed among all features, the attention is made not only to those with local correlations but also the global ones. The result is finally multiplied with the value to generate the self-attention feature maps. They are then added back to the original features to strengthen the features of interest. Such an attention mechanism is important to the training of the UR block. It allows the UR block to reconstruct images with correlated features. For instance, for the extracted features as shown in Fig. \\ref{fig:structure}, the UR block will try to reconstruct both the left and right eyes together since they have a high correlation as indicated in the training face images in the dataset. Since the correlation of the features may change at different resolutions, another UR block is used after the first one on the features of higher resolution.\n\n\\begin{figure} [htb]\n\t\\centering\n\t\\includegraphics[width=0.5\\linewidth]{Figures\/Fig_Atten_Structure.pdf}\n\t\\caption{Structure of the attention mechanism.}\n\t\\label{fig:attn}\n\\end{figure}\n\nDirect stacking covolutional and attention layers would lead to the obvious performance drop because of weak information flow and gradient vanishing \\cite{He_2016_CVPR}. Therefore, we add a shortcut connection after the convolution block and attention mechanism that forms a residual structure (shown in Fig. \\ref{fig:structure}(b)), which is widely adopted in image processing tasks \\cite{Wang_2017_CVPR, zhang2021explainable, park2021dynamic}. The residual setting avoids the gradient vanishing problem and therefore reduces the learning difficulty. As shown in Fig. \\ref{fig:structure}, the self-attention layers effectively estimate the residue between the input feature and desired one. Finally, we utilize the unshuffle layer to reshape the feature map to enlarge the spatial sizes. The effectiveness of the residual self-attention structure is further discussed in Section \\ref{Sec:ablationstudy}.\n\n\n\\subsubsection{\\textbf{Post-processing block}}\nThe Post-processing block contains two convolutional layers (shown in Fig. \\ref{fig:structure}(c)). It estimates the real $\\hat{\\mathbf{x}}_{Re}$ and imaginary $\\hat{\\mathbf{x}}_{Im}$ parts of the image from the up-sampled features. Then, the phase part $\\hat{\\mathbf{x}}$ of the image can be calculated by:\n\\begin{equation}\n\t\\label{Eq:converttophase}\n\t\\hat{\\mathbf{x}}=tan^{-1}\\left( \\frac{\\hat{\\mathbf{x}}_{Im}}{\\hat{\\mathbf{x}}_{Re}} \\right).\n\\end{equation}\n\n\n\\subsubsection{\\textbf{Loss Function}}\nWe regard the phase retrieval problem as a supervised learning task. At the training stage, we prepare a set of paired training data. Each Fourier intensity measurement has a paired complex-valued image as the ground-truth target. Then, the deep learning model can be optimized by minimizing the $\\ell_1$-norm distance between the estimation and the ground-truth images plus the TV-norm to improve the image edges. The loss function $\\mathcal{L}$ is defined as: \n\\begin{equation}\n\t\\mathcal{L} = \\left\\| \\mathbf{x}_{Re}-\\hat{\\mathbf{x}}_{Re} \\right\\|_1 + \\left\\| \\mathbf{x}_{Im}-\\hat{\\mathbf{x}}_{Im} \\right\\|_1 + \\lambda TV\\left(\\hat{\\mathbf{x}}\\right).\n\\end{equation}\nThe symbols $\\mathbf{x}_{Re}$ and $\\mathbf{x}_{Im}$ represent the ground-truth real and imaginary parts of the image, respectively. \n\n \\section{Simulation Results} \\label{Sec:simulation}\n\n\\subsection{Simulation Setup}\n\nWe conducted the simulation on two different datasets: the Real-world Affective Faces (RAF) dataset \\cite{li2017reliable} and Fashion-MNIST dateset \\cite{xiao2017\/online}. Specifically, the training data was converted to gray-scale image and resized to $128 \\times 128$. Then the pre-processed images $\\mathbf{x} \\in [0, 1]$ were used to generate complex-valued images of two forms: phase-only and magnitude-phase. For generating the phase-only images, we mapped $\\mathbf{x}$ to the $2\\pi$ phase domain via an exponential function: $\\exp(2\\pi i \\mathbf{x})$, with the magnitude kept as $1$. On the other hand, we used the images of Fashion-MNIST dataset to generate the magnitude-phase images. To be specific, we directly used the pre-processed images $\\mathbf{x} \\in [0, 1]$ as the magnitudes of the images, i.e., $\\mathbf{x}_{mag} = \\mathbf{x} \\in [0, 1]$. Then we set the phase parts through $\\mathbf{x}_{phase} = \\exp(2\\pi i \\mathbf{x}_{mag})$ and combined them by $\\mathbf{x} = \\mathbf{x}_{mag} \\circ \\mathbf{x}_{phase}$ as the magnitude-phase training and testing images, where $\\circ$ denotes the elementwise multiplication.\n\nAs mentioned in Section \\ref{Sec:Method}, we suggest to reduce the saturation problem by convolving the Fourier intensity measurement with a defocus function. Based on the convolution property of the Fourier transform, we can simulate it by multiplying the testing image with a defocus kernel $\\mathbf{h}$ in the spatial domain. It can be generated through the Fresnel propagation \\cite{goodman2017introduction}. Denote the original image as $x(p, q)$ and $X(u, v)$ as its Fourier transform, then the optical field $X_L(u', v')$ in the defocus plane is expressed as:\n\\begin{dmath}\n\t\\label{Eq:defocuskernel}\n\tX_L(u', v') \n\t= C \\iint X(u, v) H(u - \\frac{u'}{\\lambda L}, v - \\frac{v'}{\\lambda L}) du\\ dv,\n\\end{dmath}\nwhere $\\lambda$, $L$ and $H$ denotes the wavelength, distance between the lens after the object and the defocus plane, and the Fourier transform of the defocus kernel $\\mathbf{h}$, respectively. $C$ is the constant and $h(p, q)$ is given as $e^{\\frac{j\\pi}{\\lambda L}c(p^{2}+q^{2})}$. Based on the convolution theorem, Eq. \\eqref{Eq:defocuskernel} is equivalent to the elementwise multiplication between $h$ and $x$ in the spatial domain with the scaling factor $\\lambda L$. \n\nThen, oversampled discrete Fourier transform (DFT) with size $ 762\\times762 $ (the same size as in real experiments) was performed to generate the intensity measurements. For each image, the Fourier intensity measurement was capped at the value $ 4095 $ ($12$-bit) to simulate the dynamic range of the practical scientific cameras. The central $128\\times 128$ patch was cropped as the input of the \\textit{SiPRNet} (the performance of different input sizes is discussed in Section \\ref{Sec:ablainputsize}). The total epoch and mini-batch size were set as $1,000$ and $32$. The Adam optimizer \\cite{kingma2017adam} with the learning rate of $10^{-4}$ was implemented to update the gradients in each training step. Note that the learning rate was updated through \\textit{StepLR} scheduler in PyTorch every $100$ epoches with $\\gamma = 0.9$.\n\nWe adopted Peak Signal to Noise Ratio (PSNR, higher the better), Structural Similarity (SSIM, higher the better) and Mean Absolute Error (MAE, lower the better) as the performance metrics to measure the difference between the reconstructed images $\\hat{\\mathbf{x}}$ and ground truth images $\\mathbf{x}$. For each dataset, we selected the first $1,000$ images in the test set as the testing samples for evaluation. We converted the reconstructed real and imaginary parts to phase-only images via Eq. \\eqref{Eq:converttophase}. The phase values of both images were shifted by $\\pi$ to ensure non-negative values for the computation of PSNR. The average PSNR, SSIM and MAE were acquired among all $1,000$ testing images for evaluation. \n\n\\subsection{Ablation Study} \\label{Sec:ablationstudy}\n\nTo gain a deeper insight into the improvements obtained by \\textit{SiPRNet}, we conducted additional ablation studies on different components of our framework and analyzed their impact on the speed and accuracy of the network.\n\n\n\\subsubsection{\\textbf{Impact of Input Size}} \\label{Sec:ablainputsize}\n\nFor the two-dimensional Fourier spectrum, the central regions represent the low-frequency information while outer areas contain the high-frequency components. Structured signals (e.g., natural images) usually have their energy concentrated in the low-frequency regions; hence the high-frequency components have relatively lower intensity and are dominated by noise. An ablation analysis was performed to determine the size of the central regions to be used. In our experiments, the intensity measurements have the size of $762 \\times 762$. We fed the central $[32 \\times 32, 64 \\times 64, 128 \\times 128, 256 \\times 256, 512 \\times 512]$ of the intensity measurement to the same network and study the performance, as shown in Fig. \\ref{fig:Inputsize}. The average testing PSNR from $30000th$ to $35000th$ training steps was chosen for comparison. The box-plot of the performance, the floating point operations per second (FLOPs) and the numbers of trainable parameters with respect to the sizes are shown in Fig. \\ref{fig:Inputsize}. Although the model performance keeps increasing along with the increase of input size, the number of the trainable parameters also boosts quickly, which brings a huge computational burden. For example, it can be seen that the number of parameters dramatically increases from $ 19.3M $ to $ 69.6M $ when the input size changes from $128 \\times 128$ to $256 \\times 256$. While the growth speed of the performance is linear when the input size is smaller than $256 \\times 256$, it saturates as the input size increases to about $128 \\times 128$. Therefore, we adopt $128 \\times 128$ as the input size to balance performance and efficiency. \n\n\\begin{figure}[htb]\t\t\n\t\\centering\n\t\\includegraphics[width=0.7\\linewidth]{Figures\/Fig_Ablation_Input_Size.pdf}\n\t\\caption{Ablation study on the input sizes. The results are obtained with the RAF dataset \\cite{li2017reliable}.}\n\t\\label{fig:Inputsize}\n\\end{figure}\n\n\\subsubsection{\\textbf{Impact of Defocus Distance}} \\label{Sec:abladefocusdist}\n\nThe effect of the defocus distance $L$ is discussed in this section. We investigated $6$ situations of the defocus distance: $0mm$ (no defocusing), $15mm$, $20mm$, $30mm$, $45mm$, and $75mm$. The defocus kernel is a pure-phase object generated by Holoeye built-in function. The default defocus distance is $30mm$ in our proposed \\textit{SiPRNet}. The size of the input intensity measurements is $128 \\times 128$. The box-plot for average PSNR is presented in Fig. \\ref{fig:Defocusdist}. As shown in the figure, the reconstruction performance significantly degrades when there is no defocusing ($L = 0 mm$). It is because some pixels are severely saturated and give wrong information to the network. On the other hand, the performances are similar when $L = 15$ to $45 mm$ although the PSNR is slightly higher when $L = 30 mm$. It shows that the performance of the system is not sensitive to the choice of $L$ as long as it is not too small or too big. When the defocus distance is large, i.e. $75 mm$, the PSNR significantly drops. It is because a portion of the Fourier intensity measurements (high-frequency areas) become very small due to the decrease in light intensity with the long defocus distance. To balance the illumination efficiency and the reconstruction performances, the defocus distance $L$ was selected as $30 mm$ for all experiments.\n\n\\begin{figure}[htb]\t\t\n\t\\centering\t\\includegraphics[width=0.7\\linewidth]{Figures\/Fig_Ablation_DefocusDist.pdf}\n\t\\caption{Ablation study on the defocus distance. The results are obtained with the RAF dataset \\cite{li2017reliable}. }\n\t\\label{fig:Defocusdist}\n\\end{figure}\n\n\n\\subsubsection{\\textbf{Impact of Dropout Layer}} \\label{Sec:abladropout}\n\nWe investigated the effect of the number of Dropout layers in feature extraction part. We set the dropout rate as $0.2$ in every Dropout layer. Since $3$ FC layers are used in the MLP blocks, we investigated $4$ situations when introducing the Dropout layers: $0$ Dropout layer, $1$ Dropout layer after the first FC layer, $2$ Dropout layers after each of the first two FC layers, and $3$ Dropout layers after each of the FC layers. One Dropout layer is the default setting in our proposed \\textit{SiPRNet}. The size of the input intensity measurements is $128 \\times 128$. The box-plot for average PSNR is presented in Fig. \\ref{fig:Dropout}. As shown in the figure, if there is no Dropout layer, the average PSNR of the reconstructed images is the lowest because of the over-fitting problem. The network is lazy to investigate the hidden frequency representations, and the reconstruction relies on a small number of inputs. The Dropout layer randomly deactivates a part of the inputs at the training stage, which inspires the network to find more global and robust representations. At the testing stage, all input signals are used to form an ensemble learning structure that increases the generalization power of the whole network. After implementing the Dropout layer, the model performance boosts as expected. Interestingly, with more dropout layers applied, there is too much information ignored in the training stage that hurts the learning progress of the network and decreases the performance. Therefore, we adopt $30 mm$ as the defocus distance to balance performance and illumination efficiency. \n\n\\begin{figure}[htb]\t\t\n\t\\centering\t\\includegraphics[width=0.7\\linewidth]{Figures\/Fig_Ablation_Dropout.pdf}\n\t\\caption{Ablation study on the number of Dropout layers in feature extraction. The results are obtained with the RAF dataset \\cite{li2017reliable}. }\n\t\\label{fig:Dropout}\n\\end{figure}\n\n\n\\subsubsection{\\textbf{Impact of Residual Units}} \\label{Sec:ablaresidual}\n\nWe also investigated the effects of the number of residual units in the phase reconstruction part. As shown in Fig. \\ref{fig:structure}, each residual unit contains two layers of Conv Blocks interlaced with two self-attention layers plus the PReLU at the end. The size of the input images to the network is still $128 \\times 128$. With everything remaining the same but just changing the number of residual units in each UR block, the average PSNRs of the estimated images are obtained and shown in Fig. \\ref{fig:ResLayer}. It can be seen that using the residual units brings in better PSNR performance. The model with one residual unit has similar or even better performance than the models with more residual units. To pursue better PSNR and efficiency, we choose to use one residual unit in the UR block.\n\n\\begin{figure}[htb]\t\t\n\t\\centering\t\\includegraphics[width=0.7\\linewidth]{Figures\/Fig_Ablation_ResLayer.pdf}\n\t\\caption{Ablation study on the number of residual units. The results are obtained with the RAF dataset \\cite{li2017reliable}. }\n\t\\label{fig:ResLayer}\n\\end{figure}\n\n\\subsubsection{\\textbf{Model Robustness}} \\label{Sec:CrossValidation}\nBesides the generalization ability, the robustness is also a crucial metric to evaluate neural networks. The more robust the network, the less probable the overfitting problem would happen in the network. In order to demonstrate the robustness of the proposed \\textit{SiPRNet}, we utilized the K-Fold cross validation, a statistical method based on the resampling procedure \\cite{Stone_1974}. The procedure of the K-Fold cross validation can be briefly described as follows: \\textit{1.} Randomly shuffle the dataset. \\textit{2.} Split the dataset into $K$ groups. \\textit{3.} In each fold, take a group as the testing dataset and the remaining groups as the training dataset. Then, fit the deep learning model on the training dataset and do the evaluation on the testing dataset. \\textit{4.} Reset the model parameters and repeat step \\textit{3} until $K$ folds are accomplished. \\textit{5.} Summarize the model evaluation scores obtained in each fold. To be specific, we adopted $K=5$ folds for evaluation and we ran $100$ epochs in each fold. In each fold, $80\\%$ of the samples were treated as the training data and the rest $20\\%$ were utilized for testing. We used the training (for the training dataset) and testing (for the testing dataset) losses as the evaluation scores. The training loss was recorded every $20$ training batches and the testing loss was collected every epoch. The simulated results on the RAF and Fashion-MNIST datasets are shown in Fig.\\ref{Fig:KFold}. As can be seen, both losses decrease quickly with the growth of the training epoch for both datasets. The curves of the losses are similar in each fold. Although there are small variations at the beginning of the training stage ($5-15$ epochs in Fig.\\ref{Fig:KFold}(a) and $10-25$ epochs in Fig.\\ref{Fig:KFold}(b)), all folds converge to the optimal direction and thus are close to the average fold curve. In conclusion, the K-fold cross validation results demonstrate the robustness and effectiveness of \\textit{SiPRNet}. \n\n\\begin{figure}[htb]\n\t\\centering\n\t\\begin{subfigure}[t]{\\textwidth}\n\t\t\\begin{tabular}{*{2}{c@{\\extracolsep{0em}}} }\n\t\t\t\\includegraphics[width=0.5\\textwidth]{Figures\/Fig_Ablation_KFold_RAF.pdf}\n\t\t\t\\includegraphics[width=0.5\\textwidth]{Figures\/Fig_Ablation_KFold_RAF_test.pdf}\n\t\t\\end{tabular}\n\t\t\\caption{}\n\t\\end{subfigure}\n\t\n\t\\begin{subfigure}[t]{\\textwidth}\n\t\t\\begin{tabular}{*{2}{c@{\\extracolsep{0em}}} }\n\t\t\t\\includegraphics[width=0.5\\textwidth]{Figures\/Fig_Ablation_KFold_Fashion.pdf}\n\t\t\t\\includegraphics[width=0.5\\textwidth]{Figures\/Fig_Ablation_KFold_Fashion_test.pdf}\n\t\t\\end{tabular}\n\t\t\\caption{}\n\t\\end{subfigure}\n\t\n\t\\caption{Training \\& Testing errors of $5$-Fold cross validation on (a) RAF dataset, and (b) Fashion-MNIST dataset, respectively. Fold Avg. (colored in red) denotes the mean of $5$ folds.}\n\t\\label{Fig:KFold}\n\\end{figure}\n\n\n\\subsection{Simulation Results}\n\nIn this paper, we compare our \\textit{SiPRNet} with the traditional algorithms and learning-based methods. Compared methods include the traditional optimization-based algorithms such as Gerchberg\u2013Saxton (GS) \\cite{Gerchberg1972APA}, Hybrid Input-Output (HIO) \\cite{Fienup:82}, and total variation ADMM (ADMM-TV) \\cite{chang2018total}. In addition, we compare the proposed \\textit{SiPRNet} with other CNN-based methods including: Plain Network (denoted as PlainNet) \\cite{He_2016_CVPR}: the basic CNN structure that stacks convolutional layers one by one; Residual Network (denoted as ResNet) \\cite{He_2016_CVPR, Nishizaki_2020}: a popular deep-learning structure and has several skip connections among the convolutional layers; residual dense network (denoted as ResDenseNet) \\cite{Zhang_2018_CVPR}: a network that adds more skip connections to achieve better efficiency; Lensless imaging network (denoted as LenlessNet) \\cite{Sinha17}: a residual UNet to recover the phase-only images in a lensless imaging system; Conditional Generative Adversarial Phase Retrieval Network (denoted as PRCGAN) \\cite{uelwer2021phase}: a network that utilizes the conditional adversarial learning strategy for phase retrieval task; and NNPhase \\cite{Wu_cw5029}: a network to recover the complex-valued images from Fraunhofer diffraction patterns. \n\n\\textbf{\\textit{Phase-only datasets:}} The phase retrieval results with the phase-only datasets are shown in Table \\ref{table:RAFresult}, Table \\ref{table:Fashionresult}, and Fig. \\ref{fig:compvisual}. To save space, the qualitative results of PlainNet and ResDenseNet are not included in Fig. \\ref{fig:compvisual} due to their inferior performances. All the CNN-based methods were trained with the same settings on the RAF and Fashion-MNIST datasets. Note that the hyper-parameters of all traditional methods were optimally fine-tuned and fixed when the methods were evaluated with all testing images. For each traditional optimization-based method, we did three different trials of evaluation for each testing image, where each trial had $ 400 $ iterations. We manually selected the reconstructed image with the highest PSNR as the final result for comparison. For the CNN-based methods, we compared the average PSNR of reconstructed images given by the trained models at the last epoch.\n\nAs can be seen in Table \\ref{table:RAFresult} and Table \\ref{table:Fashionresult}, the proposed \\textit{SiPRNet} achieves much better performance than all compared methods. Compared with the traditional optimization-based methods, \\textit{SiPRNet} has average PSNR and SSIM gains of at least $6.544dB$ and $0.457$, respectively, on the RAF dataset (Table \\ref{table:RAFresult}). The PSNR and SSIM gains can even reach $19.58dB$ and $0.816$, respectively, on the Fashion-MNIST dataset (Table \\ref{table:Fashionresult}). This is because traditional optimization-based methods cannot converge to satisfactory results with only one intensity measurement. As shown in Fig. \\ref{fig:compvisual}, only ADMM-TV method can recover the contour of the facial images ($7$th column). In addition, \\textit{SiPRNet} can achieve PSNR and SSIM gains of at least $1.631dB$ and $0.065$, respectively, compared with other state-of-the-art CNN-based methods on the two datasets. It can be seen in Fig. \\ref{fig:compvisual} that the reconstructed images of the proposed \\textit{SiPRNet} outperform other CNN-based methods qualitatively. In addition, the error maps ($4$th column) also show that the estimated images are close to the ground-truth images. As can be seen, the reconstructed images only differ from the original ones in some local areas with minor errors, while most regions are close to $0$.\n\nNote that we train all the networks with our quantized defocused Fraunhofer diffraction patterns datasets. Although the reconstructed images via PRCGAN \\cite{uelwer2021phase} seem to have better high-frequency information (facial details), the reconstructed image does not share similarities with the original image (last second column). In fact, the performance of PRCGAN in our simulation is poorer than that reported in the original paper of PRCGAN. It is because the dataset used in the original simulation of PRCGAN does not consider the saturation problem, which is commonly found when imaging the Fourier intensity. Besides, the authors assume the object images are magnitude-only (no phase information), which is far from real situations. As for NNPhase \\cite{Wu_2021}, it aims to deal with the crystal data that the application is quite specific. Besides, LenslessNet \\cite{Sinha17} is initially designed for the phase retrieval task based on Fresnel diffraction patterns with lensless imaging systems. When testing with the more general and realistic datasets used in our simulations, the performances of these networks drop significantly compared with those reported in their original papers. We will show in the next section that these networks also suffer from the same problems in the real experiments. The above simulation results show that the proposed \\textit{SiPRNet} outperforms the traditional optimization-based and CNN-based methods both quantitatively and qualitatively. More simulation results can be found in \\nameref{Sec:Moreresults}. \n\nWe also evaluate the inference speed of different methods by collecting the mean execution time of one hundred samples. The results are presented in Table \\ref{table:RAFresult} (last column). Since our approach is designed in an end-to-end manner that can effectively utilize the huge computational power of the current advanced GPU devices, it can be seen in Table \\ref{table:RAFresult} that the proposed the \\textit{SiPRNet} is very fast and can realize real-time performance ($3.569$ ms\/img). It is comparable with other learning-based approaches. In contrast, the inference speed of the traditional algorithms is quite slow (in the order of seconds) since lots of iterations are required, and they can only be implemented with the CPU.\n\n\\textbf{\\textit{Magnitude-phase dataset:}} The simulation results are shown in Fig. \\ref{fig:ampphase}. We select HIO, NNPhase, and PRCGAN for comparison since they all perform Fourier phase retrieval. It can be seen that the proposed \\textit{SiPRNet} outperforms all compared approaches. In particular, the results show that PRCGAN cannot deal with the complex-valued phase retrieval tasks, and it is difficult for HIO to reconstruct complex-valued images with single-shot intensity measurement. The magnitude and phase images given by these approaches have poor quality. As is presented in Fig. \\ref{fig:ampphase}(b), the proposed \\textit{SiPRNet} achieves much better performance than all compared methods evaluated by different metrics. For the magnitude part, the proposed \\textit{SiPRNet} achieves average PSNR and SSIM gains of at least $2.145dB$ and $0.156$, respectively. For the phase part, the increases in PSNR and SSIM gains can reach $3.214 dB$ and $0.181$, respectively. The above simulation results show that the proposed \\textit{SiPRNet} can be generalized to complex-valued image datasets (containing both magnitude and phase parts) and outperform the existing methods.\n\n\\begin{table}[ht]\n\t\\centering\n\n\n\t\\caption{Quantitative comparison (average MAE\/PSNR\/SSIM\/Inference Time) with the state-of-the-art methods for phase retrieval on the RAF dataset. Best and second best performances are in \\textcolor{red}{red} and \\textcolor{blue}{blue} colors, respectively.}\n\t\\begin{adjustbox}{width=0.8\\columnwidth, center}\n\t\t\\begin{tabular}{|c|c|c|c|c|c|c|c|}\n\t\t\t\\hline \\multirow{2}{*}{ Methods } & \\multicolumn{2}{c|}{ MAE $\\downarrow$} & \\multicolumn{2}{c|}{ PSNR $\\uparrow$} & \\multicolumn{2}{c|}{ SSIM $\\uparrow$} & \\multirow{2}{*}{\\parbox{1.2cm}{\\centering Inference Time (ms\/img) }}\\\\\n\t\t\t\\cline { 2 - 7 } & Sim. & Exp. & Sim. & Exp. & Sim. & Exp. & \\\\\n\t\t\t\\hline\n\t\t\t\\hline GS \\cite{Gerchberg1972APA} & $0.969$ & $1.147$ & $15.099$ & $12.770$ & $0.250$ & $0.129$ & $1.233\\times 10^{3}$\\\\\n\t\t\tHIO \\cite{Fienup:82} & $1.007$ & $1.185$ & $14.854$ & $12.554$ & $0.239$ & $0.127$ & $1.226\\times 10^{3}$\\\\\n\t\t\tADMM-TV \\cite{chang2018total} & $1.119$ & $1.600$ & $12.660$ & $9.902$ & $0.206$ & $0.025$ & $1.324\\times 10^{4}$\\\\\n\t\t\tPlainNet \\cite{He_2016_CVPR} & $0.975$ & $1.180$ & $13.841$ & $12.479$ & $0.425$ & $0.377$ & $1.583$\\\\\n\t\t\tResNet \\cite{He_2016_CVPR, Nishizaki_2020} & $0.894$ & $0.812$ & $15.178$ & $15.831$ & $0.456$ & $0.471$ & $2.438$\\\\\n\t\t\tResDenseNet \\cite{Zhang_2018_CVPR} & $1.084$ & $1.010$ & $13.344$ & $13.800$ & $0.380$ & $0.387$ & $1.732$\\\\\n\t\t\tLenslessNet \\cite{Sinha17} & $\\textcolor{blue}{0.671}$ & $\\textcolor{blue}{0.704}$ & $\\textcolor{blue}{17.573}$ & $\\textcolor{blue}{17.132}$ & $\\textcolor{blue}{0.598}$ & $\\textcolor{blue}{0.596}$ & $7.396$\\\\\n\t\t\tPRCGAN \\cite{uelwer2021phase} & $0.951$ & $0.945$ & $14.738$ & $14.833$ & $0.511$ & $0.508$ & $0.983$\\\\\n\t\t\tNNPhase \\cite{Wu_cw5029} & $0.755$ & $0.931$ & $16.753$ & $15.416$ & $0.558$ & $0.544$ & $5.834$\\\\\n\t\t\t\\textbf{\\textit{SiPRNet} (Ours)} & $\\textcolor{red}{0.570}$ & $\\textcolor{red}{0.582}$ & $\\textcolor{red}{19.204}$ & $\\textcolor{red}{18.905}$ & $\\textcolor{red}{0.663}$ & $\\textcolor{red}{0.663}$ & $3.569$\\\\\n\t\t\t\\hline\n\t\t\\end{tabular}\n\t\\end{adjustbox}\n\t\\label{table:RAFresult}\n\\end{table}\n\n\\begin{table}[ht]\n\t\\centering\n\n\n\n\t\\caption{Quantitative comparison (average MAE\/PSNR\/SSIM) with the state-of-the-art methods for phase retrieval on the Fashion-MNIST dataset. Best and second best performances are in \\textcolor{red}{red} and \\textcolor{blue}{blue} colors, respectively.}\n\t\\begin{adjustbox}{width=0.7\\columnwidth,center}\n\t\t\\begin{tabular}{|c|c|c|c|c|c|c|}\n\t\t\t\\hline \\multirow{2}{*}{ Methods } & \\multicolumn{2}{c|}{ MAE $\\downarrow$} & \\multicolumn{2}{c|}{ PSNR $\\uparrow$} & \\multicolumn{2}{c|}{ SSIM $\\uparrow$} \\\\\n\t\t\t\\cline { 2 - 7 } & Sim. & Exp. & Sim. & Exp. & Sim. & Exp. \\\\\n\t\t\t\\hline\n\t\t\t\\hline GS \\cite{Gerchberg1972APA} & $1.571$ & $1.681$ & $9.852$ & $9.176$ & $0.088$ & $0.053$ \\\\\n\t\t\tHIO \\cite{Fienup:82} & $1.594$ & $1.709$ & $9.629$ & $9.082$ & $0.068$ & $0.052$ \\\\\n\t\t\tADMM-TV \\cite{chang2018total} & $0.935$ & $1.366$ & $13.772$ & $9.973$ & $0.226$ & $0.132$ \\\\\n\t\t\tPlainNet \\cite{He_2016_CVPR} & $0.589$ & $0.601$ & $16.916$ & $16.672$ & $0.629$ & $0.623$ \\\\\n\t\t\tResNet \\cite{He_2016_CVPR, Nishizaki_2020} & $0.289$ & $0.304$ & $21.394$ & $20.806$ & $0.739$ & $0.726$ \\\\\n\t\t\tResDenseNet \\cite{Zhang_2018_CVPR} & $0.381$ & $0.419$ & $19.406$ & $18.252$ & $0.675$ & $0.657$ \\\\\n\t\t\tLenslessNet \\cite{Sinha17} & $\\textcolor{blue}{0.179}$ & $\\textcolor{blue}{0.221}$ & $\\textcolor{blue}{26.784}$ & $\\textcolor{blue}{23.833}$ & $\\textcolor{blue}{0.842}$ & $\\textcolor{blue}{0.802}$\\\\\n\t\t\tPRCGAN \\cite{uelwer2021phase} & $1.598$ & $1.602$ & $10.102$ & $10.126$ & $0.422$ & $0.417$ \\\\\n\t\t\tNNPhase \\cite{Wu_cw5029} & $0.404$ & $0426$ & $23.956$ & $22.778$ & $0.739$ & $0.443$ \\\\\n\t\t\t\\textbf{\\textit{SiPRNet} (Ours)} & $\\textcolor{red}{0.127}$ & $\\textcolor{red}{0.144}$ & $\\textcolor{red}{29.209}$ & $\\textcolor{red}{28.132}$ & $\\textcolor{red}{0.884}$ & $\\textcolor{red}{0.872}$ \\\\\n\t\t\t\\hline\n\t\t\\end{tabular}\n\t\\end{adjustbox}\n\t\\label{table:Fashionresult}\n\\end{table}\n\n \n\n\\begin{figure}[htbp]\n\n\n\n\t\\begin{adjustbox}{width=1\\textwidth,center}\n\t\t\\begin{tabular}{c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}}\n\t\t\t\n\t\t\t\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/measurement_12_sim.jpg}~\n\t\t\t&\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0011_angle_gt.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0011_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/img_0011_residual.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/\/GS_RAF_simulation_12.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/\/HIO_RAF_simulation_12.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/\/ADMM_sim.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/\/Res18_simulation.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/lensless_sim_img_0011_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/PRCGAN_img_0011_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/NNPhase_img_0011_angle_sim.png}\n\t\t\t\\\\\n\t\t\t\n\t\t\t\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/measurement_12_exp.jpg}~\n\t\t\t&\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0011_angle_gt_exp.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0011_angle_exp.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/img_0011_residual_exp.png}~\t\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/GS_RAF_exp_12.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/HIO_RAF_exp_12.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/ADMM_exp.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/Res18_exp.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/lensless_exp_img_0011_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/PRCGAN_img_0011_angle_exp.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/NNPhase_img_0011_angle.png}\n\t\t\t\n\t\t\t\\\\\n\t\t\t\n\t\t\t\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/measurement_304_exp.jpg}~\n\t\t\t&\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0303_angle_gt.png}~\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0303_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/img_0303_residual.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/GS_RAF_exp_304.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/HIO_RAF_exp_304.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/ADMM_exp.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/Res18_exp2.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/lensless_exp_img_0303_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/PRCGAN_img_0303_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/NNPhase_img_0303_angle.png}\n\t\t\t\\\\\n\t\t\t\n\t\t\t\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/measurement_0847.png}\n\t\t\t&\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0847_angle_gt.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t, valign=t]{Figures\/main_result\/SiPRNet_img_0847_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t, valign=t]{Figures\/main_result\/img_0847_residual.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t, valign=t]{Figures\/main_result\/GS_img_0847_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t, valign=t]{Figures\/main_result\/HIO_img_0847_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t, valign=t]{Figures\/main_result\/ADMM_RAF_0847.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t, valign=t]{Figures\/main_result\/ResNet_img_0847_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t, valign=t]{Figures\/main_result\/LenslessNet_img_0847_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t, valign=t]{Figures\/main_result\/PRCGAN_img_0847_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t, valign=t]{Figures\/main_result\/NNPhase_img_0847_angle.png}\n\t\t\t\\\\ \t\t\t\n\t\t\t\n\t\t\t\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/measurement_646_exp_fashion.jpg}\n\t\t\t&\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0645_angle_gt.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0645_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/img_0645_residual.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/GS_Fashion_exp_646.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/HIO_Fashion_exp_646.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/ADMM_exp_fashion.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/ResNet_exp_fashion.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/lensless_exp_img_0645_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/PRCGAN_img_0645_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/NNPhase_img_0645_angle.png}\n\t\t\t\\\\\n\t\t\t\n\t\t\t\\includegraphics[height=0.09\\textwidth, valign=t]{Figures\/main_result\/measurement_0007.png}\n\t\t\t&\\includegraphics[height=0.105\\textwidth, width=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0007_angle_gt.pdf}~ &\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/SiPRNet_img_0007_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/img_0007_residual.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/GS_img_0007_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/HIO_img_0007_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/ADMM_Fashion_0007.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/ResNet_img_0007_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/LenslessNet_img_0007_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/PRCGAN_img_0007_angle.png}~\n\t\t\t&\\includegraphics[width=0.09\\textwidth, valign=t]{Figures\/main_result\/NNPhase_img_0007_angle.png}\n\t\t\t\\\\ \n\t\t\t\n\t\t\t{\\scriptsize \\makecell[c]{Fourier\\\\measurement}} & {\\scriptsize \\makecell[c]{Ground-truth \\\\ (phase)}} & {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet}}\\\\(ours)} }& {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet}} \\\\(Error Map)} } & {\\scriptsize GS } & {\\scriptsize HIO} & {\\scriptsize ADMM-TV} & {\\scriptsize ResNet}\n\t\t\t& {\\scriptsize LenslessNet} & {\\scriptsize PRCGAN} & {\\scriptsize NNPhase} \n\t\t\t\\\\\n\t\t\\end{tabular}\n\t\\end{adjustbox}\n\t\\vspace{-0.2cm}\n\t\\caption{Simulation and experimental results of different phase retrieval methods on the RAF \\cite{li2017reliable} and the Fashion-MNIST datasets \\cite{xiao2017\/online}. The first \\textbf{column} denotes the Fourier intensity measurements (pixel values: $0 \\rightarrow 4095$). The second \\textbf{column} shows the corresponding phase parts of the complex-valued images with scale bars (for experimental results) at the bottom left corners. The fourth \\textbf{column} denotes the error maps of \\textit{SiPRNet}. The other \\textbf{columns} present the reconstructed images through different methods. Except for the Fourier intensity measurements (the first \\textbf{column}), the colormap of the rest \\textbf{columns} ranges from $0$ to $2\\pi$. The first and second \\textbf{rows} are the simulation and experimental results of the same image from the RAF dataset. The third and the fourth \\textbf{rows} show the experimental performances of other images in the RAF dataset. The fifth and the sixth \\textbf{rows} indicate the experimental results of images from the Fashion-MNIST dataset. Please zoom in for better view. More results can be found in \\nameref{Sec:Moreresults}. }\n\t\\label{fig:compvisual}\n\n\\end{figure}\n\n\n\\begin{figure}[htb]\n\t\n\t\\begin{subfigure}{\\textwidth}\t\n\t\t\\centering\n\t\t\\begin{adjustbox}{width=1\\textwidth, center} \\begin{tabular}{c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}}\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/measurement_0093.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0093_magnitude_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0093_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/HIO_img_0093_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/NNPhase_img_0093_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/PRCGAN_img_0093_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0093_angle_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0093_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/HIO_img_0093_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/NNPhase_img_0093_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/PRCGAN_img_0093_angle.png}~\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/measurement_0550.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0550_magnitude_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0550_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/HIO_img_0550_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/NNPhase_img_0550_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/PRCGAN_img_0550_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0550_angle_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0550_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/HIO_img_0550_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/NNPhase_img_0550_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/PRCGAN_img_0550_angle.png}~\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/measurement_0756.png}~\n\t\t\t\t&\\includegraphics[height = 0.118\\textwidth,width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0756_magnitude_gt.pdf}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0756_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/HIO_img_0756_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/NNPhase_img_0756_magnitude.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/PRCGAN_img_0756_magnitude.png}~\n\t\t\t\t&\\includegraphics[height = 0.118\\textwidth,width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0756_angle_gt.pdf}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/SiPRNet_img_0756_angle.png}~\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/HIO_img_0756_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/NNPhase_img_0756_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/mag_phase\/PRCGAN_img_0756_angle.png}~\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t{\\scriptsize \\makecell[c]{Fourier\\\\measurement}} & {\\scriptsize \\makecell[c]{Ground-truth \\\\ (mag.)}} & {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet} } \\\\ (mag.)}} & {\\scriptsize \\makecell[c]{HIO \\\\ (mag.)}} & {\\scriptsize \\makecell[c]{NNPhase \\\\ (mag.)}} & {\\scriptsize \\makecell[c]{PRCGAN \\\\ (mag.)}} & {\\scriptsize \\makecell[c]{Ground-truth \\\\ (phase)}} & {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet} } \\\\ (phase)}} & {\\scriptsize \\makecell[c]{HIO \\\\ (phase)}} & {\\scriptsize \\makecell[c]{NNPhase \\\\ (phase)}} & {\\scriptsize \\makecell[c]{PRCGAN \\\\ (phase)}} \n\t\t\t\t\\\\\n\t\t\t\\end{tabular}\n\t\t\\end{adjustbox}\n\t\t\\vspace*{-2mm}\n\t\t\\caption{}\n\t\\end{subfigure}\n\t\n\t\\begin{subfigure}{\\textwidth}\t\n\t\t\\centering\n\t\t\\begin{adjustbox}{width=0.7\\columnwidth, center}\n\t\t\t\\begin{tabular}{|c|c|c|c|c|c|c|}\n\t\t\t\t\\hline \\multirow{2}{*}{ Methods } & \\multicolumn{2}{c|}{ MAE $\\downarrow$} & \\multicolumn{2}{c|}{ PSNR $\\uparrow$} & \\multicolumn{2}{c|}{ SSIM $\\uparrow$} \\\\\n\t\t\t\t\\cline { 2 - 7 } & Mag. & Phase & Mag. & Phase & Mag. & Phase \\\\\n\t\t\t\t\\hline\n\t\t\t\t\\hline HIO \\cite{Fienup:82} & $0.298$ & $1.898$ & $7.709$ & $8.317$ & $0.330$ & $0.043$ \\\\\n\t\t\t\tPRCGAN \\cite{uelwer2021phase} & $0.073$ & $1.654$ & $18.672$ & $9.924$ & $0.588$ & $0.387$ \\\\\n\t\t\t\tNNPhase \\cite{Wu_cw5029} & $0.109$ & $0.687$ & $18.169$ & $18.752$ & $0.395$ & $0.46$ \\\\\n\t\t\t\t\\textbf{\\textit{SiPRNet} (Ours)} & $\\mathbf{0.053}$ & $\\mathbf{0.359}$ & $\\mathbf{20.817}$ & $\\mathbf{21.966}$ & $\\mathbf{0.744}$ & $\\mathbf{0.594}$ \\\\\n\t\t\t\t\\hline\n\t\t\t\\end{tabular}\n\t\t\\end{adjustbox}\n\t\t\\caption{}\n\t\\end{subfigure}\n\t\n\t\\vspace{-0.2cm}\n\t\\caption{(a) Qualitative simulation results of different phase retrieval methods on the magnitude-phase Fashion-MNIST datasets \\cite{xiao2017\/online}. The pixel values of the Fourier intensity measurements (first column) range from $0$ to $4095$. The colormap of the magnitude parts (second to the sixth columns) ranges from $0$ to $1$, while the colormap of the phase parts (seventh to the last columns) ranges from $0$ to $2\\pi$. Please zoom in for better view. (b) Quantitative comparisons (average MAE\/PSNR\/SSIM) on the same dataset. Best performances are denoted in \\textbf{bold} font.}\n\t\\label{fig:ampphase}\n\t\n\\end{figure}\n\n \\section{Experimental Results} \\label{Sec:experiment}\n\n\\subsection{Experimental Setup}\n\nWe also evaluated our proposed \\textit{SiPRNet} with a practical optical system. A picture of the experimental setup is shown in Fig. \\ref{fig:setup}. It comprises a Thorlabs $ 10 mW$ HeNe laser with wavelength $ \\lambda = 632.8 nm$ ; and a $ 12 $-bit $ 1920\\times1200 $ Kiralux CMOS camera with pixel pitch $ 5.86 \\mu m$ . Besides, the optical system includes one $75 mm$, three $100 mm$, and one $150 mm$ lenses to form a standard imaging system. The choice of the lenses is for generating Fourier intensity images of appropriate size to match the size of the CMOS sensor. In practice, these lenses can be compactly arranged such that the whole imaging system can be made very small. However, the lens design is out of the scope of this study. A $ 1920\\times1080 $ Holoeye Pluto phase-only SLM with pixel pitch $ \\delta_{SLM}=8\\mu m $ was used to generate the testing objects in the experiments. Specifically, the images in the two datasets mentioned above were loaded to the SLM to impose multiple-level phase changes. The images were pre-multiplied with a defocus kernel; thus, there was no need to move the camera beyond the focal plane. It is to ensure that the defocusing process, which is not our main focus, will be perfectly implemented so that we only need to focus on the reconstruction performances of different methods. The defocus kernel is a pure-phase object generated by Holoeye built-in function corresponding to $L = 30 mm$. The reason for choosing $L = 30mm$ has been explained in Section \\ref{Sec:abladefocusdist}. We utilized the central $ 128\\times 128 $ SLM pixels in all experiments. The resulting Fourier intensity measurements were captured by the camera and formed the training and testing datasets for \\textit{SiPRNet}. The size of the acquired measurement is $\\frac{\\lambda L}{d}$ where $d$ denotes the pixel pitch of the SLM. \n\n\\begin{figure}[ht] \n\t\\centering\\includegraphics[width=0.8\\linewidth]{Figures\/Fig_Setup.pdf}\n\t\\caption{Experimental setup of the defocus phase retrieval system.}\n\t\\label{fig:setup}\n\t\n\\end{figure}\n\n\\subsection{Experimental Results}\n\nWith the training dataset, we fine-tuned all the pre-trained models acquired in the simulation with the experimental measurements. That is, we re-trained the models with the trained parameters from the simulation and used the experimental measurements as the new training inputs. The quantitative and qualitative experimental results are presented in Table \\ref{table:RAFresult}, Table \\ref{table:Fashionresult}, and Fig. \\ref{fig:compvisual}, respectively. \n\nIt is clearly shown in Table \\ref{table:RAFresult} and Table \\ref{table:Fashionresult} that the experimental results of \\textit{SiPRNet} only decline a little compared with the simulation results. It proves the robustness of the proposed method. As shown in Fig. \\ref{fig:compvisual}, all traditional optimization-based methods fail to reconstruct the original images because it is insufficient to converge to satisfactory solutions with only one intensity measurement for those methods. Compared with other state-of-the-art CNN-based methods, \\textit{SiPRNet} can achieve at least $1.773dB$ and $0.067$ gains in PSNR and SSIM, respectively, on the two datasets. As can be clearly seen in Fig. \\ref{fig:compvisual}, the proposed \\textit{SiPRNet} can reconstruct the most similar contours as the ground truth and preserve more high-frequency details compared with other methods. The above experimental results show that the proposed \\textit{SiPRNet} outperforms all traditional optimization-based and CNN-based methods quantitatively and qualitatively when using in practical phase retrieval systems. However, there is always room for improvement. It is shown in the results of all learning-based approaches that the reconstructed images are a bit blurry compared with the ground truths. We note that for the captured intensity images, the values of the high-frequency components are much smaller than those in the low-frequency areas. The noise in the images further makes the high-frequency components difficult to identify. It introduces much difficulty to the neural networks to extract the features of the high-frequency components. Further investigation is needed to improve the sharpness of the reconstructed images. Extra denoising and attention modules may be applied to better recover the high-frequency components.\n\n\\section{Conclusion} \n\tIn conclusion, this paper proposed an end-to-end CNN structure for single-shot maskless phase retrieval. We have shown how the optical system should be constructed to ensure the feasibility of the solution and reduce the saturation problem due to the high dynamic range of Fourier intensity measurements. We proposed a novel CNN named \\textit{SiPRNet} to reconstruct phase-only images with the central areas of the intensity measurements. To fully extract the representative features, we proposed an MLP block embedded with a Dropout layer to fuse the information of the frequency components. A self-attention structure was included to enhance the global correlation in the reconstructed phase images. The simulation and experimental results demonstrate that the proposed \\textit{SiPRNet} performs much better than the traditional phase retrieval algorithms and other CNN structures, making it a promising solution to practical phase retrieval applications. At the moment, we only conduct simulations and experiments on synthesized objects. Further research is underway to enhance the network structure for working with more complex image datasets and realistic objects. \n\n \\section*{Appendix} \\label{Sec:Moreresults}\n\n\tIn this Appendix, we show additional qualitative comparison results. Fig. \\ref{fig:firstimage}, \\ref{fig:secondimage}, \\ref{fig:thirdimage} present the simulation and experimental results. The detailed explanations are attached with each figure.\n\t\n\t \n\t\\begin{figure}[htb]\n\t\t\\begin{adjustbox}{width=1\\textwidth,center} \\begin{tabular}{c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}}\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/measurement_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/SiPRNet_img_0287_angle_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/SiPRNet_img_0287_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/img_0287_residual.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/GS_RAF_simulation_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/HIO_RAF_simulation_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/ADMM_sim_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/Res18_simulation_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/lensless_sim_img_0287_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/CGAN_simulation_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/NNPhase_img_0287_angle.png}\n\t\t\t\t\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/measurement_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/SiPRNet_img_0544_angle_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/SiPRNet_img_0544_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/img_0544_residual.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/GS_RAF_simulation_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/HIO_RAF_simulation_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/ADMM_sim_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/Res18_simulation_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/lensless_sim_img_0544_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/CGAN_simulation_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/NNPhase_img_0544_angle.png}\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/measurement_722.png}~\n\t\t\t\t&\\includegraphics[height=0.118\\textwidth,width=0.1\\textwidth, valign=t]{Figures\/first\/SiPRNet_img_0721_angle_gt.pdf}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/SiPRNet_img_0721_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/img_0721_residual.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/GS_RAF_simulation_722.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/HIO_RAF_simulation_722.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/ADMM_sim_722.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/Res18_simulation_722.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/lensless_sim_img_0721_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/CGAN_simulation_722.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/first\/NNPhase_img_0721_angle.png}\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t{\\scriptsize \\makecell[c]{Fourier\\\\measurement}} & {\\scriptsize \\makecell[c]{Ground-truth \\\\ (phase)}} & {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet}}\\\\(ours)} }& {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet}} \\\\(Error Map)} } & {\\scriptsize GS } & {\\scriptsize HIO} & {\\scriptsize ADMM-TV} & {\\scriptsize ResNet}\n\t\t\t\t& {\\scriptsize LenslessNet} & {\\scriptsize PRCGAN} & {\\scriptsize NNPhase} \n\t\t\t\t\\\\\n\t\t\t\\end{tabular}\n\t\t\\end{adjustbox}\n\t\t\\vspace{-0.2cm}\n\t\t\\caption{Simulation results of different phase retrieval methods on the RAF dataset \\cite{li2017reliable}. The first \\textbf{column} denotes the Fourier intensity measurements (pixel values: $0 \\rightarrow 4095$). The second \\textbf{column} shows the corresponding phase parts of the complex-valued images. The fourth \\textbf{column} denotes the error maps of the \\textit{SiPRNet}. The other \\textbf{columns} present the reconstructed images through different methods. Except for the Fourier intensity measurements (the first \\textbf{column}), the colormap of the rest \\textbf{columns} ranges from $0$ to $2\\pi$. Each \\textbf{row} presents the simulation results of a testing image. Please zoom in for better view.}\n\t\t\\label{fig:firstimage}\n\t\t\n\t\\end{figure}\n\t\n\t\n\t\n\n\t\\begin{figure}[htb]\n\t\t\\begin{adjustbox}{width=1\\textwidth,center} \\begin{tabular}{c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}}\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/measurement_178.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/SiPRNet_img_0177_angle_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/SiPRNet_img_0177_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/img_0177_residual.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/GS_RAF_exp_178.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/HIO_RAF_exp_178.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/ADMM_exp_178.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/Res18_exp_178.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/lensless_exp_img_0177_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/CGAN_exp_178.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/NNPhase_img_0177_angle.png}\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/measurement_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/SiPRNet_img_0287_angle_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/SiPRNet_img_0287_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/img_0287_residual.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/GS_RAF_exp_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/HIO_RAF_exp_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/ADMM_exp_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/Res18_exp_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/lensless_exp_img_0287_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/CGAN_exp_288.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/NNPhase_img_0287_angle.png}\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/measurement_545.png}~\n\t\t\t\t&\\includegraphics[height=0.118\\textwidth, width=0.1\\textwidth, valign=t]{Figures\/second\/SiPRNet_img_0544_angle_gt.pdf}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/SiPRNet_img_0544_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/img_0544_residual.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/GS_RAF_exp_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/HIO_RAF_exp_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/ADMM_exp_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/Res18_exp_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/lensless_exp_img_0544_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/CGAN_exp_545.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/second\/NNPhase_img_0544_angle.png}\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t{\\scriptsize \\makecell[c]{Fourier\\\\measurement}} & {\\scriptsize \\makecell[c]{Ground-truth \\\\ (phase)}} & {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet}}\\\\(ours)} }& {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet}} \\\\(Error Map)} } & {\\scriptsize GS } & {\\scriptsize HIO} & {\\scriptsize ADMM-TV} & {\\scriptsize ResNet}\n\t\t\t\t& {\\scriptsize LenslessNet} & {\\scriptsize PRCGAN} & {\\scriptsize NNPhase} \n\t\t\t\t\\\\\n\t\t\t\\end{tabular}\n\t\t\\end{adjustbox}\n\t\t\\vspace{-0.2cm}\n\t\t\\caption{Experimental results of different phase retrieval methods on the RAF dataset \\cite{li2017reliable}. The first \\textbf{column} denotes the Fourier intensity measurements (pixel values: $0 \\rightarrow 4095$). The second \\textbf{column} shows the corresponding phase parts of the complex-valued images with scale bars at the bottom left corners. The fourth \\textbf{column} denotes the error maps of \\textit{SiPRNet}. The other \\textbf{columns} present the reconstructed images through different methods. Except for the Fourier intensity measurements (the first \\textbf{column}), the colormap of the rest \\textbf{columns} ranges from $0$ to $2\\pi$. Each \\textbf{row} presents the experimental results of a testing image. Please zoom in for better view.}\n\t\t\\label{fig:secondimage}\n\t\t\n\t\\end{figure} \n\t\n\t\n\n\t\\begin{figure}[H]\n\t\t\\begin{adjustbox}{width=1\\textwidth, center} \\begin{tabular}{c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}c@{\\extracolsep{0em}}}\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/measurement_sim_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/SiPRNet_img_0106_angle_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/SiPRNet_img_0106_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/img_0106_residual.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/GS_Fashion_simulation_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/HIO_Fashion_simulation_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/ADMM_sim_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/ResNet_simulation_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/lensless_sim_img_0106_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/CGAN_simulation_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/NNPhase_img_0106_angle_sim.png}\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/measurement_exp_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/SiPRNet_img_0106_angle_gt_exp.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/SiPRNet_img_0106_angle_gt_exp.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/img_0106_residual_exp.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/GS_Fashion_exp_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/HIO_Fashion_exp_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/ADMM_exp_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/ResNet_exp_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/lensless_exp_img_0106_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/CGAN_exp_107.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/NNPhase_img_0106_angle.png}\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/measurement_sim_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/SiPRNet_img_0926_angle_gt.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/SiPRNet_img_0926_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/img_0926_residual.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/GS_Fashion_simulation_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/HIO_Fashion_simulation_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/ADMM_sim_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/ResNet_simulation_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/lensless_sim_img_0926_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/CGAN_simulation_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/NNPhase_img_0926_angle_sim.png}\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/measurement_exp_927.png}~\n\t\t\t\t&\\includegraphics[height=0.118\\textwidth,width=0.1\\textwidth, valign=t]{Figures\/third\/SiPRNet_img_0926_angle_gt_exp.pdf}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/SiPRNet_img_0926_angle_exp.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/img_0926_residual_exp.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/GS_Fashion_exp_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/HIO_Fashion_exp_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/ADMM_exp_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/ResNet_exp_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/lensless_exp_img_0926_angle.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/CGAN_exp_927.png}~\n\t\t\t\t&\\includegraphics[width=0.1\\textwidth, valign=t]{Figures\/third\/NNPhase_img_0926_angle.png}\n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\t{\\scriptsize \\makecell[c]{Fourier\\\\measurement}} & {\\scriptsize \\makecell[c]{Ground-truth \\\\ (phase)}} & {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet}}\\\\(ours)} }& {\\scriptsize \\makecell[c]{\\textbf{\\textit{SiPRNet}} \\\\(Error Map)} } & {\\scriptsize GS } & {\\scriptsize HIO} & {\\scriptsize ADMM-TV} & {\\scriptsize ResNet}\n\t\t\t\t& {\\scriptsize LenslessNet} & {\\scriptsize PRCGAN} & {\\scriptsize NNPhase} \n\t\t\t\t\\\\\n\t\t\t\t\n\t\t\t\\end{tabular}\n\t\t\\end{adjustbox}\n\t\t\\vspace{-0.2cm}\n\t\t\\caption{Simulation and experimental results of different phase retrieval methods on the Fashion-MNIST dataset \\cite{xiao2017\/online}. The first \\textbf{column} denotes the Fourier intensity measurements (pixel values: $0 \\rightarrow 4095$). The second \\textbf{column} shows the corresponding phase parts of the complex-valued images with scale bars (for experimental results) at the bottom left corners. The fourth \\textbf{column} denotes the error maps of \\textit{SiPRNet}. The other \\textbf{columns} present the reconstruction images through different methods. Except for the Fourier intensity measurements (the first \\textbf{column}), the colormap of the rest \\textbf{columns} ranges from $0$ to $2\\pi$. The first and second \\textbf{rows} present the simulation and experimental results of the same image. The third and the fourth \\textbf{rows} show the simulation and experimental results of another image. Please zoom in for better view.}\n\t\t\\label{fig:thirdimage}\n\t\t\n\t\\end{figure} \n\n \\section*{Data availability} The datasets used in this paper are available in \\cite{li2017reliable, xiao2017\/online}. Data underlying the results presented in this paper are not publicly available at this time but may be obtained from the authors upon reasonable request.\n \t\n \t\\section*{Funding} General Research Fund (PolyU 15225321), Centre for Advances in Reliability and Safety (CAiRS)\n \n \\section*{Acknowledgments} This work was supported by the Hong Kong Research Grant Council under General Research Fund no. PolyU 15225321, and the Centre for Advances in Reliability and Safety (CAiRS) admitted under AIR@InnoHK Research Cluster.\n \n \\section*{Disclosures} The authors declare no conflicts of interest.\n\n\t","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nCertainly the existence of black holes was one of the most interesting predictions of General Relativity. However, the discovery that these objects, due to quantum fluctuations, emit as black bodies with temperatures dictated by the surface gravity \\cite{Bekenstein:1972tm,Bekenstein:1973ur,Hawking:1974rv,Hawking:1974sw}, a purely geometric magnitude, simultaneously exposed black holes as thermodynamic objects and showed that black holes are scenarios where geometry and thermodynamics intertwine.\n\nThe first law of the black hole thermodynamics (see for instance Ref.\\cite{Wald:1984rg}) is given by\n\\begin{equation}\\label{1aleySinPdV}\ndM=TdS+\\Omega dJ+ \\phi dQ,\n\\end{equation}\nand represents the balance of energy through the modification of the macroscopic parameters of the black hole. Here $M$ corresponds to the mass parameter and was considered the internal energy of the system, namely the (ADM) mass. Finally, in Eq.(\\ref{1aleySinPdV}) $T$ is the temperature computed as the $(4\\pi)^{-1} \\kappa$ with $\\kappa$ is the surface gravity of the black hole horizon, $S$ is the entropy, $\\Omega$ is the angular velocity, $J$ is the angular momentum, $\\phi$ is the electrostatic potential and $Q$ is the electric charge.\n\nBy comparing Eq.(\\ref{1aleySinPdV}) with the first law of thermodynamics one can notice the absence of the pressure\/volume term, namely $-pdV$, which would stand for \\textit{the macroscopic work} done by the system. In a matter of speaking, this was due to the lack of the concepts of volume and pressure for a black hole in the original derivation, based on accretion processes. Several works have studied this issue in the last 20 years \\cite{Padmanabhan:2002sha,Tian:2010gn,Cvetic:2010jb,Dolan:2011xt}. To address this problem let us consider the first law of thermodynamics (see for example \\cite{Dolan:2011xt}):\n\\begin{equation}\\label{1aley1}\ndU = TdS - pdV +\\Omega dJ+ \\phi dQ,\n\\end{equation}\nwhere $U$ stands for the internal energy. To extend Eq.(\\ref{1aleySinPdV}) to match Eq.(\\ref{1aley1}), one can think of assigning a volume to the black hole by considering the \\textit{volume} defined its radius, for example $V=\\frac{4}{3}\\pi r_+^3$ for $d=4$ case. Unfortunately, since entropy is a function of the horizon radius as well, then Eq.(\\ref{1aley1}) would be inconsistent due to $dS$ and $dV$ would not be independent directions. To address this problem, in Ref.\\cite{Kastor1}, was proposed to reinterprete the mass parameter as the Enthalpy of the black hole, instead of internal energy $U$. Moreover, the cosmological constant was connected with the thermodynamic pressure. The law obtained is called the first law of (black hole) thermodynamics in the {\\it extended phase space}. In Ref.\\cite{MotorTermico1}, on the other hand, the promotion of the mass parameter to the enthalpy is based on the fact that to form a black hole would require to cut off a region of the space, and therefore an initial energy equal to $E_0 = -\\rho V$, with $\\rho$ the energy density of the system, is needed. In four dimensions the thermodynamics volume corresponds to $V=\\frac{4}{3}\\pi r_+^3$. Moreover, the presence of the cosmological constant defines $\\rho = -P$ and therefore the mass parameter can be considered equivalent to\n\\begin{equation}\nM=U-\\rho V = U+pV,\n\\end{equation}\nHere $M$ is to be recognized as the enthalpy $H$ of the system. With this in mind, the first law, in this extended phase space, yields\n\\begin{equation} \\label{1aley2}\ndM=dH=TdS+Vdp+\\phi dQ+\\Omega dJ.\n\\end{equation}\nThis extended first law (\\ref{1aley2}) also was derived in reference \\cite{Mann1} by using Hamiltonian formalism. It is worth to mention that the definition of the extended phase space have allowed to construct a {\\it heat engine} in terms of a black hole, see some examples in references \\cite{MotorTermico1,MotorTermico5,MotorTermico7,MotorTermico13,MotorTermico14,MotorTermico15}, adding a new layer to our understanding of the black hole thermodynamics. Another interesting applications is the Joule Thompson expansion for black holes studied in \\cite{JT1,JT2,JT3,JT4}.\n\n\\subsection{Phase Transitions}\nThe study of phase transitions in black hole physics has called a renewed attention in the last years due to the AdS\/CFT conjecture. For instance, it is well known that the {\\it Hawking-Page phase transition} \\cite{Hawking:1982dh}, in the context of the AdS\/CFT correspondence, has been re-interpreted as the plasma gluon confinement\/deconfinement phase transition in the would-be dual (conformal) field theory. Similarly, the AdS Reissner Nordstr\\\"om's transitions in the $(\\phi-q)$ diagram have been interpreted as liquid\/gas phase transitions of Van der Waals fluids \\cite{Chamblin:1999tk,Chamblin:1999hg}.\n\nRecently the analysis of the $p-V$ critical behaviors (in the extended phase space) have been under studied extensively\\cite{RegularBHquimica, PV1,PV2,PV3,PV4,PV5,PV6,PV7,PV8,PV9,PV10,PV11,PV12,Hendi:2017fxp,Hendi:2015cqz,Hendi:2015hoa,Hendi:2014kha}. For instance, in \\cite{Mann2} was studied in the context of the charged $4D$ AdS black holes how the phase transitions between small\/large black hole are analogous to liquid\/gas transitions in a Van der Waals fluid. Moreover, it was also shown that the critical exponents, near the critical points, recovers those of Van der Waals fluid with the same \\textit{compressibility factor} $Z=3\/8$. In reference \\cite{Mann1} was introduced a \\textit{new interpretation} of the Hawking Page phase transition \\cite{Hawking:1982dh} mentioned above, but in the context of $p-V$ critical behavior. \n\n\\subsection{Higher Dimensions, Lovelock and thermodynamics}\n\nDuring the last years 50 years several branches of theoretical physics have noticed that considering higher dimensions is plausible. Now, considering higher dimension in gravity opens up a range of new possibilities that retain the core of the Einstein gravity in four dimensions. Lovelock is a one of these possibilities as, although includes higher powers of curvature corrections, its equations of motion are of second order and thus causality is still insured. Generic Lovelock theory is the sum of the Euler densities $L_n$ in lower dimensions ($2n \\leq d$) multiplied by coupling constants $\\alpha_n$ \\footnote{In even dimensions the maximum order, $n=d\/2$, is in turn the corresponding Euler density of the dimension and therefore does not contribute to the equations of motion,}. The Lagrangian is\n\\begin{equation}\n \\displaystyle L = \\sum_{n=0}^{[d\/2]} \\alpha_n L_n,\n\\end{equation}\nwhere $L_n=\\frac{1}{2^{n}}\\,\\sqrt{-g}\\ \\delta_{\\nu_{1}...\\nu_{2n}}^{\\mu_{1}...\\mu_{2n}}\\ R_{\\mu_{1}\\mu_{2}}^{\\nu_{1}\\nu_{2}}\\cdots R_{\\mu_{2n-1}\\mu_{2n}}^{\\nu_{2n-1}\\nu_{2n}}$. The $L _0 \\propto 1$ term is related by the cosmological constant, $L_1 \\propto R $ is related by the Ricci scalar, and $L_2 \\propto R^{\\alpha\\beta}_{\\hspace{2ex}\\mu\\nu}R^{\\mu\\nu}_{\\hspace{2ex}\\alpha\\beta}-4 R^{\\alpha\\nu}_{\\hspace{2ex}\\beta\\nu} R^{\\beta\\mu}_{\\hspace{2ex}\\alpha\\mu} + R^{\\alpha\\beta}_{\\hspace{2ex}\\alpha\\beta}R^{\\mu\\nu}_{\\hspace{2ex}\\mu\\nu}$ is the Gauss Bonnet density. The EOM are:\n\\begin{align}\\label{LovelockGenerico}\n\\frac{1}{\\sqrt{g}}\\frac{\\delta }{\\delta g_{\\mu\\nu}} (L\\sqrt{g}) &= G^{\\mu}_{ \\hspace{1ex} \\nu} \\nonumber \\\\\n&= \\sum_{n} \\alpha_{n} \\frac{1}{2^n} \\delta^{\\mu_1 \\ldots \\mu_{2n} \\mu}_{\\nu_1 \\ldots \\nu_{2n} \\nu} R^{\\nu_{1} \\nu_{2}}_{\\hspace{2ex}\\mu_{1} \\mu_{2}}\\ldots R^{\\nu_{2n-1} \\nu_{2n}}_{\\hspace{2ex}\\mu_{2n-1} \\mu_{2n}}\n\\end{align}\nwhere $G^{\\mu \\nu}_{(LL)}$ which, by definition, satisfies $\\nabla_\\mu G^{\\mu\\nu}_{(LL)} \\equiv 0$.\n\nIn higher curvature the $p-V$ criticality have been studied for example in reference \\cite{Mann3} for generic Lovelock black holes with conformal scalar hair arising from coupling of a real scalar field to the dimensionally extended Euler densities and in reference \\cite{Nam:2019clw} for a regular charged AdS black hole in Einstein Gauss Bonnet gravity. See also \\cite{Hendi:2015soe,Xu:2014tja}.\n\nOne drawback of a generic Lovelock gravity is the existence of more than single ground state, namely more than a single constant curvature spaces solution, or equivalently more than a single potential \\textit{effective cosmological constants} \\cite{Camanho:2011rj}. However, there are two families of Lovelock gravities that have indeed a single ground state. The first family has been originally studied in \\cite{Banados:1994ur} and has a unique but $k$ degenerated ground state. \n\nThe second case is usually called {\\it Pure Lovelock gravity}. In this case the Lagrangian is a just the $n$-single term of Lagrangian plus a cosmological constant, \\textit{i.e.}, $L = \\alpha_n L_n + \\alpha_0 L_0$. This has a single and non-degenerated ground state of negative curvature $-1\/l^2$ for $n$ odd. Several solutions of Pure Lovelock theories are known. For instance vacuum black hole solutions can be found in references \\cite{Cai:2006pq,DadhichBH,Dadhich3,PureLovelock1,PureLovelock2,Aranguiz:2015voa} and regular black holes in references \\cite{milko1,milko2}. See other applications in references \\cite{Dadhich1,Dadhich2,Dadhich4,Dadhich5}.\n\nIn the next sections the thermodynamic of the Pure Lovelock solutions will be carried out in details. For this, the cosmological constant will be promoted to the intensive thermal pressure of the system. Under this assumption, it will be shown that the system behaves as Van der Waals fluid. Finally, the critical coefficient near the critical points will be computed.\n\n\\subsubsection{{\\bf Vacuum black hole in Pure Lovelock gravity: Uncharged Asymptotically AdS Case}}\n\nLet us consider a static geometry, represented in Schwarzschild coordinates by\n\\begin{equation}\\label{linelement}\nds^2 =-f(r) dt^2+ \\frac{dr^2}{f(r)} + r^2 d \\Omega^2_{d-2}.\n\\end{equation}\nFor simplicity the cosmological constant is fixed such that\n\\begin{equation} \\label{Lambda1}\n \\Lambda = - \\frac{(d-1)(d-2)}{2 l^{2n}}.\n\\end{equation}\nThis is similar to the definition in \\cite{Aranguiz:2015voa} but a numerical factor. The equations of motion reduce to merely\n\\begin{equation}\n\\frac{d}{dr} \\left ( r^{d-2n-1}\\big (1-f(r) \\big)^n + \\dfrac{ r^{d-1}}{{l}^{2n}} \\right ) =0.\n\\end{equation}\nSolutions of these equations have been studied in references \\cite{Cai:2006pq,Dadhich1}. In terms of Eq.(\\ref{linelement}) the solution is defined by\n\\begin{equation} \\label{funcionPLvacio}\nf(r) =1 - \\left (\\frac{2M}{r^{d-2n-1}} - \\frac{r^{2n}}{l^{2n}} \\right )^{1\/n} .\n\\end{equation}\nOne can notice that $f(r)$ can take complex values for even $n$. Because of this, in this work will be considered only odd $n$. The thermodynamic pressure can be read off this definition as\n\\begin{equation} \\label{p1}\n p=-\\frac{\\Lambda}{(d-2)\\Omega_{d-2}},\n\\end{equation}\nwhere $\\Omega_{d-2}$ is the unitary area of a $d-2$ sphere. It is worth to stress that Eqs.(\\ref{Lambda1},\\ref{p1}) coincide with the definitions in Ref.\\cite{Mann1,Mann2} for $n=1$ and $d=4$. Now, replacing Eqs.(\\ref{Lambda1},\\ref{p1}) into Eq.(\\ref{funcionPLvacio}) yields\n\\begin{equation} \\label{funcionPLvacio1}\n f(r)=1 - \\left (\\frac{2M}{r^{d-2n-1}} - \\frac{2 \\Omega_{d-2}r^{2n}}{d-1}p \\right )^{1\/n},\n\\end{equation}\nwhere the dependence of $f(r)$ on $p,M$ have been made explicit.\n\n\\subsubsection{{\\bf Charged Pure Lovelock solution}}\nThe charged case is slightly difference since it is necessary to solve the Maxwell equations and to include an energy momentum tensor into the gravitational equations. Let us start by defining $A_\\mu = A_t(r) \\delta_\\mu^t$, which defines only the non-vanishing component of the Maxwell tensor \n\\begin{equation} \\label{PotencialVectorial}\n F_{tr}=-\\partial_rA_t(r), \n\\end{equation}\nand the single non vanishing component of the Maxwell equations, \n\\begin{equation} \\label{TensorElectromagnetico}\n \\nabla_{\\mu}F^{\\mu\\nu} =0 \\rightarrow \\frac{d}{dr} \\left ( r^{d-2} F_{tr} \\right) =0.\n\\end{equation}\nLikewise, the gravitational equations are reduced to the single independent component\n\\begin{equation}\\label{GravitationalEMEq}\n \\frac{d}{dr} \\left ( r^{d-2n-1}\\big (1-f(r) \\big)^n + \\dfrac{ r^{d-1}}{{l}^{2n}} \\right ) = (F_{tr})^2 r^{d-2}.\n\\end{equation}\nIt is straightforward to integrate Eqs.(\\ref{TensorElectromagnetico},\\ref{GravitationalEMEq}) yielding \\cite{Aranguiz:2015voa}\n\\begin{equation} \\label{funcionPLvacioQ}\nf(r) =1 - \\left (\\frac{2M}{r^{d-2n-1}} - \\frac{r^{2n}}{l^{2n}} - \\frac{Q^2}{(d-3) r^{2d-2n-4}} \\right )^{1\/n},\n\\end{equation}\nand\n\\begin{equation}\n F_{tr}=\\frac{Q}{r^{d-2}} \\rightarrow A_t(r)= \\phi_\\infty + \\frac{Q}{(d-3)r^{d-3}},\n\\end{equation}\nAs usual $\\phi_{\\infty}=0$ is fixed such that $\\displaystyle \\lim_{r\\rightarrow \\infty} A_t(r) =0$. \n\nNow, by direct observation, one can notice that for even $n$ the existence of certain ranges of $r$ where $f(r)$ can take complex values. To avoid this only odd $n$ will be considered from now on.\n\nAs done previously, by replacing Eqs(\\ref{Lambda1},\\ref{p1}) into Eq.(\\ref{funcionPLvacioQ}), $f(r)$ can be written in terms of the thermodynamics variables as\n\\begin{equation} \\label{funcionPLvacio1Q}\n f(r)=1 - \\left (\\frac{2M}{r^{d-2n-1}} - \\frac{2 \\Omega_{d-2}r^{2n}}{d-1}p - \\frac{Q^2}{(d-3) r^{2d-2n-4}} \\right )^{1\/n}.\n\\end{equation}\n \n\n\\section{Extended phase space in Vacuum Pure Lovelock Gravity}\n\nIn this section the extended phase space formalism, mentioned above, will be introduced to analyze the solutions. One can notice that the fist law of the thermodynamics, Eq.(\\ref{1aley2}), for the non rotating case takes the form\n\\begin{equation} \\label{1aleyconcarga}\n dH= dM=TdS+Vdp + \\phi dQ.\n\\end{equation}\nNow, in order to construct a thermodynamic interpretation one must notice that under any transformation of the parameters the function $f(r_+,M,Q,l)$ must still vanishes, otherwise the transformation would not be mapping black holes into black holes in the space of solutions. Indeed, $\\delta f(r_+,M,Q,l) = 0$ and $f(r_+,M,Q,l)=0$ are to be understood as constraints on the evolution along the space of parameters. However, there is another approach by recalling that the mass parameter, $M$, is also to be understood as a function of the parameters $M(r_+,l,Q)$ as well.\n\nSince the thermodynamic parameters are $S,p$ and $Q$, therefore it is convenient to reshape $M=M(S,p,Q)$ in order to explicitly obtain \n\\begin{equation} \\label{1aley3}\n dM= \\left ( \\frac{\\partial M}{\\partial S} \\right)_{p,Q}dS + \\left ( \\frac{\\partial M}{\\partial p} \\right)_{S,Q}dp + \\left ( \\frac{\\partial M}{\\partial Q} \\right)_{p,S}dQ .\n\\end{equation}\nThis corresponds to the definitions of the component of the tangent vector in the space of parameters, but also correspond to the definitions of the temperature, thermodynamic volume and electric potential in the form of\n\\begin{eqnarray}\nT &=&\\left ( \\frac{\\partial M}{\\partial S} \\right)_{p,Q},\\label{temperatura} \\\\\nV &=& \\left ( \\frac{\\partial M}{\\partial p} \\right)_{S,Q} \\textrm{ and }\\label{volumen} \\\\\n\\phi &=& \\left ( \\frac{\\partial M}{\\partial Q} \\right)_{S,p}. \\label{potencial}\n\\end{eqnarray}\n\nOn the other hand, the variation along the space of parameters of the condition defined by $f(r_+,M,p,Q)=0$,\n\\begin{eqnarray}\n df(r_+,M,p,Q) &=& 0\\label{variacion}\\\\\n &=& \\frac{\\partial f}{\\partial r_+} dr_+ + \\frac{\\partial f}{\\partial M} dM + \\frac{\\partial f}{\\partial p} dp + \\frac{\\partial f}{\\partial Q} dQ,\\nonumber\n\\end{eqnarray}\nyields a second expression for $dM$ given by\n\\begin{align} \\label{variacion1}\n dM &= \\left (\\frac{1}{4\\pi} \\frac{\\partial f}{\\partial r_+} \\right ) \\left( -\\frac{1}{4 \\pi} \\frac{\\partial f}{\\partial M} \\right)^{-1} dr_+ \\nonumber \\\\\n &+ \\left( - \\frac{\\partial f}{\\partial M} \\right)^{-1} \\left ( \\frac{\\partial f}{\\partial p} \\right ) dp +\\left( - \\frac{\\partial f}{\\partial M} \\right)^{-1} \\left ( \\frac{\\partial f}{\\partial Q} \\right ) dQ,\n\\end{align}\nwhich must coincide with equation (\\ref{1aley3}). In Eq.(\\ref{variacion1}) one can recognize presence of the temperature, which geometrically is defined as\n\\begin{equation} \\label{temperaturaconocida}\n T= \\frac{1}{4\\pi} \\frac{\\partial f}{\\partial r_+},\n\\end{equation}\n yielding \n\\begin{equation} \\label{dS1}\n\\left( -\\frac{1}{4 \\pi} \\frac{\\partial f}{\\partial M} \\right)^{-1} dr_+= dS,\n\\end{equation}\nwhich is a very known result. Now, by the same token, the thermodynamic volume and electric potential are given by\n\\begin{equation} \\label{volumen1}\n V=\\left ( \\frac{\\partial M}{\\partial p} \\right)_{S,Q}= \\left( - \\frac{\\partial f}{\\partial M} \\right)^{-1} \\left ( \\frac{\\partial f}{\\partial p} \\right ) \n\\end{equation}\nand\n\\begin{equation} \\label{potencial1}\n \\phi = \\left ( \\frac{\\partial M}{\\partial Q} \\right)_{S,p}= \\left( - \\frac{\\partial f}{\\partial M} \\right)^{-1} \\left ( \\frac{\\partial f}{\\partial Q} \\right ),\n\\end{equation}\nrespectively. These expressions will be discussed in the next sections.\n\n\\subsection{Smarr Expression}\nConsidering Plank units, one can notice that $p$, the pressure, has units of $[p]=\\ell^{-2n}$. See Eqs.(\\ref{Lambda1},\\ref{p1}, \\ref{presionPL1}). Likewise, one can check that $[M] = \\ell^{d-2n-1}$, $[Q] = \\ell^{d-n-2}$ and $[S] = \\ell^{d-2n}$. \n\nFollowing {\\it Euler's theorem} \\cite{Altamirano:2014tva}, with $M(S,p,Q)$, one can construct the {\\it Smarr formula} for Pure Lovelock gravity given by\n\\begin{equation}\n \\frac{d-2n-1}{d-2n}M=TS+\\frac{d-n-2}{d-2n} \\phi Q - \\frac{2n}{d-2n}VP.\n\\end{equation}\nThis coincides with the definitions discussed in \\cite{Altamirano:2014tva,Mann2,Mann1} for $n=1$ . Derivations of the Smarr formula for generic Lovelock theory are discussed in \\cite{Dolan:2014vba,Kastor:2010gq}. \n\n\\subsection{ Uncharged asymptotically AdS}\nReplacing the solution of Eq.(\\ref{funcionPLvacio1}) into Eq.(\\ref{dS1}) yields\n\\begin{align}\n\\left( -\\frac{1}{4 \\pi} \\frac{\\partial f}{\\partial M} \\right)^{-1} dr_+&=2\\pi n r_+^{d-2n-1} \\nonumber \\\\\n&=d \\left ( \\frac{2}{d-2n}n \\pi r_+^{d-2n} \\right ) = dS,\n\\end{align}\nor equivalently\n\\begin{equation} \\label{entropia1}\n S= \\frac{2}{d-2n}n \\pi r_+^{d-2n},\n\\end{equation}\nwhich coincides with reference \\cite{Cai:2006pq}. Eq.(\\ref{volumen1}), on the other hand, yields\n\\begin{equation} \\label{volumen11}\n \\left( - \\frac{\\partial f}{\\partial M} \\right)^{-1} \\left ( \\frac{\\partial f}{\\partial p} \\right ) = V = \\frac{\\Omega_{d-2}}{d-1}r_+^{d-1},\n\\end{equation}\nwhich corresponds to the volume of a $(d-2)$ sphere of radius $r_+$ . This coincides with the definition in Refs.\\cite{Mann1,Mann2}.\n\n\\subsubsection{{ \\bf Fluid equation of state}}\nIt is worth to notice at this point that the temperature, see Ref.\\cite{milko1}, can be expressed as \n\\begin{equation} \\label{TemperaturaPL1}\n 4\\pi n T= \\frac{d-2n-1}{r_+} + (d-1)\\frac{ r_+^{2n-1}}{l^{2n}},\n\\end{equation}\nwhere, the temperature has units of $\\ell^{-1}$. One can make explicit the dependence on the pressure, by inserting Eqs.(\\ref{Lambda1},\\ref{p1}) into Eq.(\\ref{TemperaturaPL1}). Therefore, \n\\begin{equation} \\label{presionPL1}\n 2 \\Omega_{d-2}p = 4 \\pi n \\frac{T}{r_+^{2n-1}} - \\frac{d-2n-1}{r_+^{2n}}.\n\\end{equation}\n\n\\subsubsection{\\bf Physical Pressure}\nBefore to proceed a digression is necessary. As mentioned above $[p] = \\ell^{-2n}$, however the physical pressure, $p_G$, must be satisfied [Force\/Area]$=\\ell_p^{-2}$, since the area has units of [area]$=\\ell_p^{d-2}$. In the literature $p_g$ is called the geometrical pressure. Since $p$ and $p_G$ must be connected by $p_G \\sim \\alpha_n p$, namely $[p_{G}] = \\left [\\alpha_n p \\right]=\\ell_p^{-d}$, therefore\n \\begin{itemize}\n \\item the {\\it coupling constants} must satisfy $[\\alpha_n]=\\ell^{2n-d}$ \\cite{Oliva1}. For $n=1$ this coincides with the inverse of the higher dimensional Newton constant $G_d^{-1}$ \\cite{Maartens:2003tw} and\n \\item there is still room for a dimensionless constant which can be used to adjust the definition. \n\\end{itemize}\nWith this in mind, $p_G$ can be taken as \n\\begin{equation}\n p_{G} =2 \\Omega_{d-2} \\ell_p^{2n-d} p = 4 \\pi n \\frac{T}{\\ell_p^{d-2n}r_+^{2n-1}} + \\ldots,\n\\end{equation}\nwhich, in turns, defines the \\textit{specific volume} of the system as\n\\begin{equation} \\label{volumenespecificoPL}\n v=\\frac{\\Omega_{d-2}}{2\\pi} \\ell_p^{d-2n}r_+^{2n-1}.\n\\end{equation}\n\nNotice that $v$ indeed has units of volume, namely $[v]= \\ell_p^{d-1}$ \\cite{Hennigar:2018cnh}. Finally, after some replacements, \n\\begin{equation} \\label{presionPL2}\n p = \\frac{n T}{v} - \\frac{1}{2\\Omega_{d-2}(2\\pi\/\\Omega_{d-2})^{2n\/(2n-1)}}\\frac{d-2n-1}{v^{2n\/(2n-1)}},\n\\end{equation}\nwhich can be recognized as the Van der Walls equation for $n=1$, namely $P=T\/(v-b)-a\/v^2$, with $b=0$ \\cite{Mann1}. \n\n\\subsubsection{{\\bf Hawking-Page phase transition}}\n\nIn reference \\cite{Hawking:1982dh} was analyzed the thermodynamic behavior of Schwarzschild AdS space, which differs from the Schwarzschild case due to the presence of the AdS gravitational potential, namely the presence of $\\sim r^2\/l^2$ in $f(r)$ \\cite{Czinner:2015eyk}. In this case there is a phase transition between black hole and AdS radiation at a critical temperature $T_{HP}$ where the Gibbs free energy, $G=M-TS$ vanishes \\cite{Mann2,Czinner:2015eyk,Wang:2019vgz}. \n\nFig.\\ref{HP} displays the numerical behavior of the Gibbs free energy v\/s temperature for $n=3$ and $d=10$ for Pure Lovelock. The upper curve represents the unstable small black hole (namely with negative heat capacity) and the lower curve represents the stable large black hole. It is direct to show that this behavior is similar for any other set of values of $n$ and $d$. \n\nThe generic behavior is displayed in Fig.\\ref{HPgenerico}. The upper green curve represents the unstable small black hole, the lower blue curve represents the stable large black hole, and the orange curve represents the thermal AdS radiation whose Gibbs free energy vanishes \\cite{Mann1,Czinner:2015eyk,Wang:2019vgz,Costa:2015gol}. One can notice that for $[T_{min},T_{HP}[$ the large stable black hole has positive Gibbs free energy, therefore, the preferred state corresponds to the thermal AdS radiation. On the other hand, for $T>T_{HP}$, the preferred state is the large stable black hole whose Gibbs free energy is negative. This hints the existence of a Hawking Page phase transition between radiation and the large black hole states at $T=T_{HP}$. \n\nFinally, it is worth to notice, from Fig.\\ref{HP}, that the value of $T_{HP}$ increases as the pressure increases. This behavior is similar to the HP phase transition for the Schwarzschild AdS black holes \\cite{Mann1,Belhaj:2015hha} or the polarized AdS black holes \\cite{Costa:2015gol}.\n\n\\begin{figure}\n\\centering\n{\\includegraphics[width=3in]{HP1.eps}}\n\\caption{HP phase transitions for $n=3$, $d=10$ and $Q=1$, $p=0,0000045$(blue), $p=0,0000050$(red), $p=0,0000055$ (green). T (horizontal axis) v\/s G (vertical axis). }\n\\label{HP}\n\\end{figure}\n\n\\begin{figure}\n\\centering\n{\\includegraphics[width=3in]{HPgenerico.eps}}\n\\caption{Generic behavior of Gibbs free energy (vertical axis) v\/s T (horizontal axis). HP phase transition at $T_{HP}$}\n\\label{HPgenerico}\n\\end{figure}\n\n\\subsection{ Charged Pure Lovelock solution}\nBy replacing Eq.(\\ref{funcionPLvacio1Q}) into Eqs.(\\ref{dS1},\\ref{volumen1}) the expressions for the entropy (\\ref{entropia1}) and the volume (\\ref{volumen11}) are obtained . On the other hand, replacing Eq.(\\ref{funcionPLvacio1Q}) into Eq.(\\ref{potencial1}) yields \n\\begin{equation} \\label{potencial2}\n \\phi= \\frac{Q}{(d-3)r^{d-3}}.\n\\end{equation}\nThis thermodynamics quantities, obtained by mean of variation of parameters, are consistent with the values presented in the literature. \n\n\\subsubsection{\\bf Fluid equation of state}\nIn this case the temperature can be written as\n\\begin{equation} \\label{TemperaturaPL1Q}\n 4\\pi n T= \\frac{d-2n-1}{r_+} + (d-1)\\frac{ r_+^{2n-1}}{l^{2n}}- \\frac{Q^2}{r_+^{2d-2n-3}} .\n\\end{equation}\nBy inserting equations (\\ref{Lambda1},\\ref{p1}) into equation (\\ref{TemperaturaPL1Q}) is obtained\n\\begin{equation} \\label{presionPL2Q}\n 2 \\Omega_{d-2}p = 4 \\pi n \\frac{T}{r_+^{2n-1}} - \\frac{d-2n-1}{r_+^{2n}} + \\frac{Q^2}{r_+^{2d-4}}.\n\\end{equation}\nNow, by using the definition of the specific volume, defined in Eq.(\\ref{volumenespecificoPL}) with $\\ell_p=1$,\n\\begin{eqnarray}\n p &=& \\frac{n T}{v} - \\frac{1}{2\\Omega_{d-2}(2\\pi\/\\Omega_{d-2})^{2n\/(2n-1)}}\\frac{d-2n-1}{v^{2n\/(2n-1)}} \\nonumber \\\\\n &+& \\frac{1}{2\\Omega_{d-2}(2\\pi\/\\Omega_{d-2})^{(2d-4)\/(2n-1)}}\\frac{Q^2}{v^{(2d-4)\/(2n-1)}}. \\label{presionPL3}\n\\end{eqnarray}\nOne can notice the similarity of this relation with the corresponding one for a Van der Waals fluid for $n=1$ but with the presence of an addition third term. \n\n\\subsubsection{\\bf Critical points and compressibility factor}\nTo compare Eq.(\\ref{presionPL3}) with the behavior of a Van der Waals fluid it is necessary to determine the critical points of the system. The second order critical points are defined by the conditions, \n\\begin{equation}\n \\frac{\\partial p}{\\partial v}=0 \\textrm{ and } \\frac{\\partial^2 p}{\\partial v^2}=0.\n\\end{equation}\nThese determine the critical values\n\\begin{equation}\n v_c = \\frac{\\Omega_{d-2}}{2\\pi}\\left ( \\frac{2d^2-2dn-7d+4n+6}{n(d-2n-1)}Q^2 \\right )^{(2n-1)\/(2d-2n-4)} ,\n\\end{equation}\nand\n\\begin{align}\n T_c &= \\frac{n(d-2n-1)}{2\\pi n (2n-1)(2\\pi\/\\Omega_{d-2})^{1\/(2n-1)}v^{1\/(2n-1)}} \\nonumber \\\\\n &-\\frac{(d-2)Q^2(2\\pi\/\\Omega_{d-2})^{(2n-2d+4)\/(2n-1)}v^{(2n-2d+4)\/(2n-1)}}{2\\pi n (2n-1)(2\\pi\/\\Omega_{d-2})^{1\/(2n-1)}v^{1\/(2n-1)}}\n\\end{align}\nNotice that $p_c=p(v_c,T_c)$ can be determined from equation (\\ref{presionPL3}) by evaluation on the critical values $T_c$ and $v_c$. \n\nIn the table (\\ref{tabla1}) critical values $v_c$, $T_c$ and $p_c$ and the {\\it compressibility factor} $Z=p_cv_v\/T_c$ are displayed for different values of $n$ and $d$ . For $n=1$ and $d=4$ the compressibility factor has the exact value $Z=3\/8$ which coincides with the value of the compressibility for a Van der Waals fluid. In general, compressibility factor can be written as\n\\begin{equation} \\label{factorcompresibilidad}\n Z=\\frac{2d-2n-3}{4(d-2)},\n\\end{equation}\nwhich implies that $Z<1$ extrictly. This implies that this can be interpreted as a low pressure gas. \n\\begin{table*} \n\\caption{Critical values and Compressibility Factor ($Z$) .}\n\\begin{tabular*}{\\textwidth}{@{\\extracolsep{\\fill}}lccccc@{}}\n\\hline\n\\multicolumn{1}{c}{$n$ } & \\multicolumn{1}{c}{$d$ } & \\multicolumn{1}{c}{$v_c$}& \\multicolumn{1}{c}{$T_c$}&\\multicolumn{1}{c}{$p_c$}&\\multicolumn{1}{c}{$Z=\\dfrac{p_cv_c}{T_c}$} \\\\\n \\hline\n$1$ & $4$ & $\\Omega_{d-2}\/(2\\pi)Q\\sqrt{6}$ & $1\/(18\\pi Q) \\sqrt{6}$ &$1\/(24\\Omega_{d-2}Q^2)$ & $3\/8$ \\\\\n\\hline \n$1$ & $5$ & $\\Omega_{d-2}\/(4\\pi)(120Q^2)^{1\/4}$ & $4\/(75\\pi)(6750\/Q^2)^{1\/4}$ &$2\/(45\\Omega_{d-2}Q)\\sqrt{30}$ & $5\/12$ \\\\\n\\hline\n$1$ & $6$ & $\\Omega_{d-2}\/(6\\pi)(6804Q^2)^{1\/6}$ & $9\/(98\\pi)(201.684\/Q^2)^{1\/6}$ &$9\/(112\\Omega_{d-2})(294\/Q^2)^{1\/3}$ & $7\/16$ \\\\\n\\hline\n$3$ & $8$ & $\\Omega_{d-2}\/(2\\pi)(14Q^2)^{5\/6}$ & $3\/(490\\pi) (14^5\/Q^2)^{1\/6}$ &$1\/(280\\Omega_{d-2}Q^2)$ & $7\/24$ \\\\\n\\hline\n$3$ & $9$ & $\\Omega_{d-2}\/(4\\pi)(21^5\\cdot 2^3Q^{10})^{1\/8}$ & $8\/(945\\pi) (21^7\\cdot 2\/Q^2)^{1\/8}$ &$4\/(735\\Omega_{d-2})(21\\cdot 2^3\/Q^6)^{1\/4}$ & $9\/28$ \\\\\n\\hline\n$3$ & $10$ & $\\Omega_{d-2}\/(18\\pi)(792)^{1\/2}Q$ & $3\/(242\\pi) (2^7\\cdot 11^9 \\cdot 9\/Q^2)^{1\/10}$ &$9\/(704\\Omega_{d-2})(726\/Q^6)^{1\/5}$ & $11\/32$ \\\\\n\\hline\n$5$ & $12$ & $\\Omega_{d-2}\/(2\\pi)(22Q^2)^{9\/10}$ & $5\/(2178\\pi)(22^9\/Q^2)^{1\/10}$ &$1\/(792\\Omega_{d-2}Q^2)$ & $11\/40$ \\\\\n\\hline\n\\end{tabular*} \\label{tabla1}\n\\end{table*}\n\n\\subsubsection{\\bf $P-v$ curve.}\n\nIn Fig.\\ref{p-v} is displayed for $n=3$ and $d=10$ the behavior of the curve $p-v$ defined by equation (\\ref{presionPL3}). Although this is an example still this behavior is generic for any values of $n$ and $d$. \n\nFor values of temperature $T>T_{c}$, the second and third factors of Eq.(\\ref{presionPL3}) are negligible in comparison with the first one, and therefore the curve approximates the form $p \\cdot v \\propto T$ and thus mimicking the behavior of an ideal gas. Conversely, for $Tv_{max}$ with $p(v_1)=p(v_2)$. In fluid theory these two solutions are known as the {\\it Van der Waals loop} and physically this corresponds to a {\\it vapor liquid equilibrium}, highlighting that phase transitions take place. \n\n\\begin{figure}\n\\centering\n{\\includegraphics[width=3in]{p-v.eps}}\n\\caption{$p-v$ curve for $T_1p_{c}$, in figure \\ref{FigTemperatura}(b) for $p=p_{c}$ and in figure \\ref{FigTemperatura}(c) for $pp_{c}$ the temperature is an increasing function of $r_+$. For $p=p_{c}$ the temperature has one inflexion point at $r_+=r_{infl}$. More relevant for this discussion is the case for $pp_{c}$ .]{\\includegraphics[width=75mm]{TemperaturaMayor.eps}} \n\\subfigure[Temperature for $p=p_{c} \\approx 0,001608$ ]{\\includegraphics[width=75mm]{TemperaturaCritica.eps}}\n\\subfigure[Temperature for $p=0,00001p_c$, in figure \\ref{CalorEspecificoMayoroigual}(b) for $p=p_c$ and in figure \\ref{CalorEspecificoMenor} for $pp_c$ the heat capacity is a positive increase function of $r_+$, and there is no phase transition, thus black hole is always stable.\n \\item For $p=p_c$ small and large black hole coexist at the inflexion point $r_+=r_{infl}$, where the heat capacity $C \\to \\infty$.\n \\item For the $p0$ defining a small stable black hole. Next, there is small unstable ($C<0$) region for $r \\in ]r_{max},r_{min}[$. Finally there is a third region $r_+>r_{min}$ where the system is a large stable black hole $(C>0)$. This hints the existence of phase transitions but, by means of the following analysis of the Gibbs free energy, one can check that only the \\textit{small stable bh\/large stable bh} transition is allowed. \n\\end{itemize}\n\n\\begin{figure}\n\\centering\n\\subfigure[Heat capacity for $p=1>p_c$.]{\\includegraphics[width=75mm]{Cmayor.eps}}\n\\subfigure[Heat capacity for $p=p_c \\approx 0,001608 $]{\\includegraphics[width=75mm]{Ccritica.eps}}\n\\caption{Behavior of Heat Capacity for $n=3$ and $d=10$ with $Q=1$ .}\n\\label{CalorEspecificoMayoroigual}\n\\end{figure}\n\n\n\\begin{figure}\n\\centering\n\\subfigure[Heat capacity for $r_+ \\in (0.896,1.6)$.]{\\includegraphics[width=75mm]{CMenor1.eps}}\n\\subfigure[Heat capacity for $r_+ \\in (2.7,3.5)$]{\\includegraphics[width=75mm]{CMenor2.eps}}\n\\caption{Behavior of Heat Capacity for $n=3$ and $d=10$ with $Q=1$, $p=0,00001p_c$, in figure \\ref{GibbsT}(b) for $p=p_c$ and in \\ref{GibbsT}(c) for $pp_{c}$ .]{\\includegraphics[width=63mm]{GibbsMayor1.eps}} \\hspace{15mm}\n\\subfigure[ Gibbs for $p=p_{c} \\approx 0,001608$ ]{\\includegraphics[width=63mm]{GibbsCritico1.eps}}\n\\subfigure[Gibbs for $p=0,00001T_c$. For $T k$ and for each $L\\subset V$ with $|L|0$ strongly biconnected components. Let $(u,w)$ be an edge in $E\\setminus E_s$ such that $u,w $ are in not in the same strongly\n\tbiconnected component of $G_s$. Then the directed subgraph $(V,E\\cup \\left\\lbrace (u,w) \\right\\rbrace )$ contains at most $t-1$ strongly biconnected components.\n\\end{lemma}\n\\begin{proof}\n\tSince $G_s$ is strongly connected, there exists a simple path $p$ from $w$ to $u$ in $G_s$. Path $p$ and edge $(u,w)$ form a simple cycle. Consequently, the vertices $u,w$ are in the same strongly biconnected component of the subgraph $(V,E\\cup \\left\\lbrace (u,w) \\right\\rbrace )$.\n\\end{proof}\n\n\\begin{figure}[htbp]\n\t\\begin{myalgorithm}\\label{algo:approximationalgorithmfor2sb}\\rm\\quad\\\\[-5ex]\n\t\t\\begin{tabbing}\n\t\t\t\\quad\\quad\\=\\quad\\=\\quad\\=\\quad\\=\\quad\\=\\quad\\=\\quad\\=\\quad\\=\\quad\\=\\kill\n\t\t\t\\textbf{Input:} A $2$-vertex strongly biconnected directed graph $G=(V,E)$ \\\\\n\t\t\t\n\t\t\t\\textbf{Output:} a $2$-vertex strongly biconnected subgraph $G_{2s}=(V,E_{2s})$\\\\\n\t\t\t{\\small 1}\\> find a minimal $2$-vertex-connected subgraph $G_{1}=(V,E_{1})$ of $G$.\\\\\n\t\t\t{\\small 2}\\> \\textbf{if} $G_{1}$ is $2$-vertex strongly biconnected \\textbf{then} \\\\\n\t\t\t{\\small 3}\\>\\>output $G_{1}$\\\\\n\t\t\t{\\small 4}\\> \\textbf{else}\\\\\n\t\t\t{\\small 5}\\>\\>$E_{2s} \\leftarrow E_{1}$\\\\\n\t\t\t{\\small 6}\\>\\> $G_{2s}\\leftarrow (V,E_{2s})$ \\\\\n\t\t\t{\\small 7}\\>\\> identify the b-articulation points of $G_{1}$.\\\\\n\t\t\t{\\small 8}\\>\\> \\textbf{for} evry b-articulation point $b\\in V$ \\textbf{do} \\\\\n\t\t\t{\\small 9}\\>\\>\\> \\textbf{while} $G_{2s}\\setminus\\left\\lbrace b \\right\\rbrace $ is not strongly biconnected \\textbf{do}\\\\ \n\t\t{\\small 10}\\>\\>\\>\\> calculate the strongly biconnected components of $G_{2s}\\setminus\\left\\lbrace b \\right\\rbrace $\\\\\n\t\t{\\small 11}\\>\\>\\>\\> find an edge $(u,w) \\in E\\setminus E_{2s}$ such that $u,w $ are not in \\\\\n\t\t{\\small 12}\\>\\>\\>\\>the same strongly biconnected components of $G_{2s}\\setminus\\left\\lbrace b \\right\\rbrace $.\\\\\n\t\t{\\small 13}\\>\\>\\>\\> $E_{2s} \\leftarrow E_{2s} \\cup\\left\\lbrace (u,w)\\right\\rbrace $ \\\\\n\t\t{\\small 14}\\>\\>\\>output $G_{2s}$\n\t\n\t\t\\end{tabbing}\n\t\\end{myalgorithm}\n\\end{figure}\n\n\\begin{lemma} \nAlgorithm \\ref{algo:approximationalgorithmfor2sb} returns a $2$-vertex strongly biconnected directed subgraph.\n\\end{lemma}\t\n\\begin{proof}\nIt follows from Lemma \\ref{def:addingedgestoreducesbcs}.\n\\end{proof}\n\nThe following lemma shows that each optimal solution for the M22VSBSS problem has at least $2n$ edges.\n\\begin{lemma} \\label{def:optsolutionatlesttwon}\n\tLet $G=(V,E)$ be a $2$-vertex-strongly biconnected directed graph. Let $O\\subseteq E$ be an optimal solution for the M2VSBSS problem. Then $|O|\\geq 2n$. \n\t\n\\end{lemma}\n\\begin{proof}\n\tfor any vertex $x\\in V$, the removal of $x$ from the subgraph $(V,O)$ leaves a strongly biconnected directed subgraph. Since each strongly biconnected directed graph is stronly connected,\n\tthe subgraph $(V,O)$ has no strong articulation points. Therefore, the directed subgraph $(V,O)$ is $2$-vertex-connected.\n\\end{proof}\nLet $l$ be the number of b-articulation points in $G_{1}$. The following lemma shows that Algorithm \\ref{algo:approximationalgorithmfor2sb} has an approximation factor of $(2+l\/2)$.\n\\begin{theorem}\n\tLet $l$ be the number of b-articulation points in $G_{1}$. Then, $|E_{2s}|\\leq l(n-1)+4n$.\n\\end{theorem}\n\\begin{proof}\n\tResults of Edmonds \\cite{Edmonds72} and Mader \\cite{Mader85} imply that $|E_1|\\leq 4n$ \\cite{CT00,Georgiadis11}. Moreover, by Lemma \\ref{def:optsolutionatlesttwon}, every optimal solution for the M22VSBSS problem has size at least $2n$. For every b-articulation point in line $8$, Algorithm \\ref{algo:approximationalgorithmfor2sb} adds at most $n-1$ edge to $E_{2s}$ in while loop. Therefore, $|E_{2s}|\\leq l(n-1)+4n$\n\\end{proof}\n\\begin{Theorem}\n\tThe running time of Algorithm \\ref{algo:approximationalgorithmfor2sb} is $O(n^{2}m)$.\n\\end{Theorem}\n\\begin{proof}\n\tA minimal $2$-vertex-connected subgraph can be found in time $O(n^2)$ \\cite{Georgiadis11,GIK20}. B-articulation points can be computed in $O(nm)$ time.\n\tThe strongly biconnected components of a directed graph can be identified in linear time \\cite{WG2010}. \n\tFurthermore, by Lemma \\ref{def:addingedgestoreducesbcs}, lines $9$--$13$ take $O(nm)$ time.\n\\end{proof}\n\n\\section{Open Problems}\n Results of Mader \\cite{Mader71,Mader72} imply that the number of edges in each minimal $k$-vertex-connected undirected graph is at most $kn$ \\cite{CT00}. Results of Edmonds \\cite{Edmonds72} and Mader \\cite{Mader85} imply that the number of edges in each minimal $k$-vertex-connected directed graph is at most $2kn$ \\cite{CT00}. These results imply a $2$-approximation algorithm \\cite{CT00} for minimum $k$-vertex-connected spanning subgraph problem for undirected and directed graphs \\cite{CT00} because every vertex in a $k$-vertex-connected undirected graphs has degree at least $k$ and every vertex in a $k$-vertex-connected directed graph has outdegree at least $k$ \\cite{CT00}. Note that these results imply a $7\/2$ approximation algorithm for the M2VSBSS problem by calculating a minimal $2$-vertex-connected directed subgraph of a $2$-vertex strongly biconnected directed graph $G=(V,E)$ and a minimal $3$-vertex connected undirected subgraph of the underlying graph of $G$.\n \n \\begin{lemma}\nLet $G=(V,E)$ be a $2$-vertex strongly biconnected directed graph. Let $G_1=(V,L)$ be a minimal $2$-vertex-connected subgraph of $G$ and let $G_2=(V,U)$ be a minimal $3$-vertex-connected subgraph of the underlying graph of $G$. Then the directed subgraph $G_s=(V,L\\cup A )$ is $2$-vertex strongly connected, where $A=\\left\\lbrace (v,w) \\in E \\text{ and } (v,w) \\in U \\right\\rbrace $. Moreover, $|L\\cup A|\\leq 7n$\n \\end{lemma}\n\\begin{proof}\n\tLet $w$ be any vertex of the subgraph $G_s$. Since the $G_1=(V,L)$ is $2$-vertex-connected, subgraph $G_s$ has no strong articulation points. Therefore, $G_s\\setminus\\left\\lbrace w\\right\\rbrace $ is strongly connected. Moreover, the underlying graph of $G_s\\setminus\\left\\lbrace w\\right\\rbrace $ is biconnected because the underlying graph of $G_s$ is $3$-vertex-connected. Results of Edmonds \\cite{Edmonds72} and Mader \\cite{Mader85} imply that $|L|<4n$. Results of Mader \\cite{Mader71,Mader72} imply that $|U|\\leq 3n$.\n\\end{proof}\n \n An open problem is whether each minimal $2$-vertex strongly biconnected directed graph has at most $4n$ edges.\n\n Cheriyan and Thurimella \\cite{CT00} presented a $(1+1\/k)$-approximation algorithm for the minimum $k$-vertex-connected spanning subgraph problem for directed and undirected graphs. The algorithm of Cheriyan and Thurimella \\cite{CT00} has an approximation factor of $3\/2$ for the minimum $2$-vertex-connected directed subgraph problem. Let $G=(V,E)$ be a $2$-vertex strongly biconnected directed graph and let $E^{CT}$ be the output of the algorithm of Cheriyan and Thurimella \\cite{CT00}. The directed subgraph $(V,E^{CT})$ is not necessarily $2$-vertex strongly biconnected. But a $2$-vertex strongly biconnected subgraph can be obtained by performing the following third phase. For each edge $e\\in E\\setminus E^{CT}$, if the underlying graph of $G\\setminus \\left\\lbrace e\\right\\rbrace $ is $3$-vertex-connected, delete $e$ from $G$. We leave as open problem whether this algorithm has an approximation factor of $3\/2$ for the M2VSBSS problem.\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\chapter{Acknowledgement}\\label{Acknowledgement}\n\n\nIt seems incredible but the Ph.D. really comes to the end! \n\nFor sure this has been the most intense period of my life both by the scientific point of view, with {\\em my} Di-Bosoni, the Higgs, tons of new things learned and discovered, and also by the personal point of view,\nwith thousands of miles across the World and several houses and cities where I spent my time. \n\nThe real risk was to get lost along the way but, luckily, I have been continuously surrounded, helped and supported by so so many nice people that it is really hard to thank everybody in one page.\n\nFor sure I couldn't had made it without the support of Prof. Giorgio Chiarelli (who now knows me by the far 2007); the work wouldn't had been possible without him but neither without the continuous exchange of ideas with Dr. Sandra Leone and within all the CDF Pisa Group.\n\nFermilab has been the other indispensable ingredient: the best place in the Chicago suburbs to do Physics and also all the rest. I met there a great research team (all the WHAM'ers), wonderful scientists and researchers from whom I learned so much.\n\nHowever there have not just been Physics in these years\\footnote{Although often friendship, study, work, spare time and Physics all overlaps in a indistinguishable way...}, my family has been always with me also from the other side of the ocean or farther away. No possibility to get anywhere without their help. Then there is Celine, ``la Piccolina'', that I met once, lost, and luckily I met again and now we are together!\n\nAnd finally there are all the friends from {\\em Sancasciani}, {\\em Cottonwood}, {\\em Siena}, {\\em Chicago}, {\\em Pisa}, the {\\em old} and the {\\em new Summer Students}, the {\\em ControTV} group and others I am surely missing. I spend with all of you a very good time and first or later I will meet again everybody scattered all around the World.\n\nThis travel lasted three years (well almost four), a pretty long period in the lifespan of somebody that is not even thirty but it flowed rapidly like a mountain torrent.\n\n\n\n\\chapter[First Evidence of Diboson in $\\ell\\nu+ HF$ ($7.5~fb^{-1}$)]{First Evidence of Diboson in Lepton plus Heavy Flavor Channel ($7.5~\\textrm{fb}^{-1}$)}\\label{App:7.5}\n\nThe first evidence of diboson production in the $\\ell \\nu + HF$ final state was already obtained in the Summer of 2011 with a preliminary version of this analysis performed on a smaller dataset, $7.5$~fb$^{-1}$ of CDF data. \n\nThe differences with respect to the work described in the main parts of the thesis are the followings:\n\\begin{itemize}\n\\item a dataset of $7.5$~fb$^{-1}$ of CDF data.\n\\item The forward electrons category (PHX) was not used.\n\\item A previous version~\\cite{CHEP_svm} of the SVM multi-jet rejection algorithm was employed. The algorithm was optimized only for binary classification and it was not possible to use the shape of the output distribution. Therefore the \\mbox{${\\not}{E_T}$} distribution was used to in the normalization of the $W + $ jets and the multi-jet backgrounds.\n\\item No flavor separator ($KIT-NN$) information was used for single-tag events. Therefore only the di-jet invariant mass distribution was used as the final signal to background discriminator both for the single and double tagged signal regions.\n\\item The statistical analysis of the significance of the observation was performed in a different way. The likelihood ratio~\\cite{mclimit} of the signal and {\\em test} hypothesis, after the fit over the nuisance parameter, was used\\footnote{This method can not be easily applied to the significance estimate of two signals (for example $WW$ {\\em vs} $WZ\/ZZ$) therefore, in the main part of the thesis, we moved to the method described in Section~\\ref{sec:sigma_eval} for the statistical analysis}. \n\\item We also estimated $95$\\% Confidence Level (CL) limits, both in the case of diboson signal and no diboson signal. The first case can be used to constrain new physics models which produces an increase of the TGC couplings.\n\\end{itemize}\nA summary of the event selection, the background estimate and the statistical analysis is reported in the following.\n\n\n\n\\section{Event Selection and Background Estimate}\\label{sec:Data}\n\nWe select events consistent with the $\\ell\\nu+HF$ signature. \n\nThe charged lepton candidate online and offline identification is described in Chapter~\\ref{chap:sel} but we consider only the tight central lepton candidates (CEM, CMUP, CMX), the loose lepton candidates (BMU, CMU, CMP, CMIO, SCMIO, CMXNT) and the isolated track candidates (ISOTRK). Loose leptons and ISOTRK are classified together in the EMC category. The $W\\to \\ell \\nu$ selection is completed by a cut on the \\mbox{${\\not}{E_T}$} variable, corrected for the presence of muons and jets. We require a \\mbox{\\mbox{${\\not}{E_T}$} $>20$~GeV} for CEM and EMC leptons while we relax the cut down to \\mbox{\\mbox{${\\not}{E_T}$} $>10$~GeV} for the tight muons categories. \n\nThe $HF$ jets selection is also identical to the one reported in Chapter~\\ref{chap:sel}: two central ($|\\eta_{Det}| < 2.0$) jets with E$_T^{cor} > 20$~GeV (energy corrected for detector effects) \non which we require the identification of one or two secondary vertices with the \\texttt{SecVtx} $b$-tagging algorithm. The $b$-tagging requirement divides the selected sample in a {\\em pretag} control region (no $b$-tag requirement) and two signal regions characterized by exactly one or two $b$-tagged jets.\n\nA relevant difference in the event selection comes from the multi-jet background rejection strategy. We used a previous version of the SVM algorithm~\\cite{CHEP_svm} based on the following six variables (described in Section~\\ref{sec:var_sel}):\n\\begin{itemize}\n\\item Lepton $p_T$, \\mbox{${\\not}{E_T}$} , $MetSig$, $\\Delta\\phi(lep,\\mbox{${\\not}{E_T^{raw}}$})$, $E_T^{raw,jet2}$ and $E_T^{cor,jet2}$.\n\\end{itemize}\nThe discriminant was optimized to work in the central region of the detector and, as binary classifier, the performances were similar to the present algorithm described in Appendix~\\ref{chap:AppSvm}.\n\nThe background estimate was done with the same methodology described in Chapter~\\ref{chap:bkg}. The only difference is that the fits, used to estimate $W + $ jets and multi-jet normalizations in the different lepton categories and tag regions, were performed on the \\mbox{${\\not}{E_T}$} distributions. Figure~\\ref{fig:qcd_fit_met} shows the result of the pretag maximum likelihood fit on the \\mbox{${\\not}{E_T}$} , pretag sample, distribution of the different lepton categories: CEM, CMUP, CMX, EMC.\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{App7.5\/Figs\/PretafFit_CEM}\n\\includegraphics[width=0.495\\textwidth]{App7.5\/Figs\/PretafFit_CMUP}\n\\includegraphics[width=0.495\\textwidth]{App7.5\/Figs\/PretafFit_CMX}\n\\includegraphics[width=0.495\\textwidth]{App7.5\/Figs\/PretafFit_ISOTRK}\n\\caption[\\mbox{${\\not}{E_T}$} Fit Used for $W +$ Jets Estimates in the Pretag Sample]{$W + $ jets and non-$W$ fraction estimates, on the pretag sample, with a maximum likelihood fit on the \\mbox{${\\not}{E_T}$} distributions. The non-$W$ background is shown in pink while the $W + $ jets component is in green. The dashed line is the sum of all the backgrounds and the points are the data. The figures show (left to right and top to bottom) the CEM, CMUP, CMX and EMC \n charged lepton categories.} \\label{fig:qcd_fit_met}\n\\end{center}\n\\end{figure}\n\nThe total background estimate is reported in Tables~\\ref{tbl:1tag} and~\\ref{tbl:2tags}. The final statistical analysis of the selected events is performed on the $M_{Inv}(jet1, jet2)$ distribution of four channels: single and double tagged candidates for the central tight leptons (CEM+CMUP+CMX) and EMC leptons. Figure~\\ref{fig:mjj_1tag} shows the high-statistics single-tag $M_{Inv}(jet1, jet2)$ distribution of central tight leptons and of the EMC leptons.\n\n\\begin{table}[htb]\n \\begin{center}\n \\begin{tabular}{lcccc}\n \\toprule\n \\multicolumn{5}{c}{\\bf Single-tag Event Selection}\\\\\n Lepton ID & CEM & CMUP & CMX & EMC \\\\ \n \\midrule \n\n Pretag Data & 61596 & 29036 & 18878 & 27946 \\\\\\midrule\n $Z+$jets &27.9 $\\pm$ 3.5 & 43.0 $\\pm$ 5.5 & 27.3 $\\pm$ 4.2 & 65.0 $\\pm$ 8.4 \\\\ \n $t\\bar{t}$ &201.3 $\\pm$ 19.6 &109.8 $\\pm$ 10.7 & 55.0 $\\pm$ 7.1 & 171.9 $\\pm$ 16.9 \\\\ \n Single Top $s$ &52.9 $\\pm$ 4.8 &28.2 $\\pm$ 2.6 & 14.0 $\\pm$ 1.8 & 38.0 $\\pm$ 3.5 \\\\ \n Single Top $t$ &71.4 $\\pm$ 8.4 &37.4 $\\pm$ 4.4 & 19.8 $\\pm$ 2.9 & 49.5 $\\pm$ 5.8 \\\\ \n $WW$ &68.0 $\\pm$ 9.4 &33.3 $\\pm$ 4.6 & 20.3 $\\pm$ 3.3 & 38.4 $\\pm$ 5.3 \\\\ \n $WZ$ &21.8 $\\pm$ 2.3 & 11.5 $\\pm$ 1.25 & 7.4 $\\pm$ 1.0 & 14.1 $\\pm$ 1.6 \\\\ \n $ZZ$ &0.44 $\\pm$ 0.04 & 0.65 $\\pm$ 0.06 & 0.42 $\\pm$ 0.05 & 0.86 $\\pm$ 0.08 \\\\ \n $W + b\\bar{b}$ &632.9 $\\pm$ 254.2 & 309.8$\\pm$ 124.1 & 192.3 $\\pm$ 77.1 & 308.9 $\\pm$ 124.3 \\\\\n $W + c\\bar{c}$ &331.0 $\\pm$ 133.7 & 155.1 $\\pm$ 62.5 & 96.2 $\\pm$ 38.8 & 164.2 $\\pm$ 66.4 \\\\\n $W + cj$ &259.9 $\\pm$ 105.0 & 127.8 $\\pm$ 51.5 & 75.3 $\\pm$ 30.4 & 106.4 $\\pm$ 43.0 \\\\\n $W+LF$ &605.2 $\\pm$ 71.3 & 283.8 $\\pm$ 31.7 & 181.0 $\\pm$ 20.6 & 346.2 $\\pm$ 39.2 \\\\ \n Non-$W$ &173.9 $\\pm$ 69.6 & 45.8 $\\pm$ 18.3 & 2.8 $\\pm$ 1.1 & 100.9 $\\pm$ 40.4 \\\\ \n \\midrule \n {\\bf Prediction} &{\\bf 2446 $\\pm$ 503} & {\\bf 1186 $\\pm$ 242 } &{\\bf 692 $\\pm$ 149} & {\\bf 1404 $\\pm$ 243} \\\\\n {\\bf Observed} &{\\bf 2332} & {\\bf 1137} & {\\bf 699} & {\\bf 1318} \\\\\n \\midrule \n {\\bf Dibosons} &{\\bf 89.7 $\\pm$ 10.2} & {\\bf 44.8 $\\pm$ 5.05} & {\\bf 27.7 $\\pm$ 3.9} & {\\bf 52.5 $\\pm$ 5.9} \\\\\n \\bottomrule \n \\end{tabular}\n\n \\caption[Observed and Expected Events with One \\texttt{SecVtx} Tags]{Summary of observed and expected events with one \\texttt{SecVtx} tag, in the $W+2$ jets sample, in \n 7.5 fb$^{-1}$ of data.}\\label{tbl:1tag}\n \\end{center}\n\\end{table}\n\n\\begin{table}[htb]\n\\begin{center}\n \\begin{tabular}{lcccc}\n \\toprule\n \\multicolumn{5}{c}{\\bf Double-tag Event Selection}\\\\\n Lepton ID & CEM & CMUP & CMX & EMC \\\\ \n \\midrule \n \n Pretag Data & 61596 & 29036 & 18878 & 27946 \\\\\\midrule\n $Z+$jets & 0.9 $\\pm$ 0.1 &2.0 $\\pm$ 0.3 & 1.2 $\\pm$ 0.2 & 3.1 $\\pm$ 0.4 \\\\ \n $t\\bar{t}$ & 42.2 $\\pm$ 6.1 &22.2 $\\pm$ 3.2 & 11.1 $\\pm$ 1.9 & 34.4 $\\pm$ 5.0 \\\\\n Single Top $s$ & 14.1 $\\pm$ 2.0 &7.6 $\\pm$ 1.1 & 3.7 $\\pm$ 0.6 & 10.2 $\\pm$ 1.4 \\\\ \n Single Top $t$ & 4.2 $\\pm$ 0.7 &2.3 $\\pm$ 0.4 & 1.2 $\\pm$ 0.2 & 3.1 $\\pm$ 0.5 \\\\ \n $WW$ & 0.6 $\\pm$ 0.1 & 0.26 $\\pm$ 0.07 & 0.16 $\\pm$ 0.04 & 0.33 $\\pm$ 0.08\\\\ \n $WZ$ & 4.0 $\\pm$ 0.6 & 1.9 $\\pm$ 0.3 & 1.4 $\\pm$ 0.2 & 2.4 $\\pm$ 0.4 \\\\ \n $ZZ$ & 0.06 $\\pm$ 0.01 & 0.12 $\\pm$ 0.02 & 0.09 $\\pm$ 0.01 & 0.16 $\\pm$ 0.02 \\\\\n $W + b\\bar{b}$ & 81.9 $\\pm$ 33.2 &42.2 $\\pm$ 17.1 & 23.4 $\\pm$ 9.5 & 44.9 $\\pm$ 18.2\\\\\n $W + c\\bar{c}$ & 4.7 $\\pm$ 1.9 &2.3 $\\pm$ 1.0 & 1.3 $\\pm$ 0.5 & 2.8 $\\pm$ 1.1 \\\\ \n $W + cj$ & 3.7 $\\pm$ 1.5 &1.9 $\\pm$ 0.8 & 1.0 $\\pm$ 0.4 & 1.8 $\\pm$ 0.7 \\\\ \n $W+LF$ & 3.2 $\\pm$ 0.7 &1.6 $\\pm$ 0.3 & 0.9 $\\pm$ 0.2 & 2.2 $\\pm$ 0.4 \\\\ \n Non-$W$ & 7.9 $\\pm$ 3.2 &4.8 $\\pm$ 1.9 & 0.1 $\\pm$ 0.5 & 0.0 $\\pm$ 0.5 \\\\ \n \\midrule \n {\\bf Prediction} & {\\bf 167.3 $\\pm$ 38.0} &{\\bf 88.9 $\\pm$ 19.6} & {\\bf 45.4 $\\pm$ 10.9} & {\\bf 105.3 $\\pm$ 21.5}\\\\\n {\\bf Observed} & {\\bf 147} &{\\bf 74} &{\\bf 39} & {\\bf 106} \\\\ \n \\midrule\n {\\bf Dibosons} &{\\bf 4.6 $\\pm$ 0.6} &{\\bf 2.1 $\\pm$ 0.3} & {\\bf 1.5 $\\pm$ 0.2} &{\\bf 2.7 $\\pm$ 0.4} \\\\\n \\bottomrule\n \\end{tabular} \n\n \\caption[Observed and Expected Events with Two \\texttt{SecVtx} Tags]{Summary of observed and expected events with two \\texttt{SecVtx} tags, in the $W+2$ jets sample, in $7.5$~fb$^{-1}$ of data.}\\label{tbl:2tags}\n\\end{center}\n\\end{table}\n\n\n\\begin{figure} [!ht]\n\\centering\n\\includegraphics[width=0.495\\textwidth]{App7.5\/Figs\/1svt_cen}\n\\includegraphics[width=0.495\\textwidth]{App7.5\/Figs\/1svt_emc} \n\\caption[$M_{Inv}(jet1,jet2)$ Distribution for Single \\texttt{SecVtx} Tag Events]{$M_{Inv}(jet1,jet2)$ distribution for the single \\texttt{SecVtx} tagged events. Tight leptons (CEM+CMUP+CMX combined) on the left and EMC leptons on the right. The best fit values for the rate and shape of the backgrounds are used in the figures.} \n\\label{fig:mjj_1tag}\n\\end{figure}\n\n\n\n\\section{Statistical Analysis}\n\nAt the time of this analysis, the process $WZ\/WW \\rightarrow \\ell \\nu $ + Heavy Flavors was not yet observed at hadron colliders. Therefore we started by evaluating a 95\\% exclusion CL on a potential signal and then we evaluated the cross section of the process and its significance. \n\nMost of the statistical analysis procedure has been described in Chapter~\\ref{chap:StatRes}. We build a likelihood function (Equation~\\ref{eq:likelihood_full}) with templates derived from the selected data and backgrounds, both shape and rate systematics are taken into account as described in Section~\\ref{sec:sys_desc}. The only relevant change concerns the $b$-tag SF for $c$-marched jets where we applied the same prescriptions of the $WH$ CDF search, it does not double the uncertainty for $c$-matched quarks. \n\nThe last relevant difference is the evaluation of the 95\\% CLs for the diboson signal. The CLs are evaluated by integrating the likelihood distribution over the unknown parameter $\\alpha$ (i.e. the diboson cross section) up to cover 95\\% of the total possible outcomes. In formulas:\n\\begin{equation}\n\\int^{\\bar{\\alpha}}_{0} \\mathscr{L}(\\alpha)\\mathrm{d}\\alpha. \n\\end{equation}\nwhere $\\mathscr{L}(\\alpha)$ is derived from Equation~\\ref{eq:likelihood_full} after the integration over the nuisance parameters, $\\vec{\\alpha}$ is reduced to just one dimension because we perform CLs on only one signal and the limit on the integration, $\\bar{\\alpha}$, is given by the condition of 95\\% coverage.\n\nA first set of expected CLs are obtained assuming no SM diboson production and generating Pseudo Experiments (PEs) on the base of the expected background yields varied within the assigned systematics. \nCombining single--tagged and double--tagged results for all lepton categories, we find an expected limit of:\n\\begin{equation}\n\\left( 0.575^{+0.33}_{-0.31} \\right) \\times \\textrm{SM prediction.}\n\\end{equation}\n\nWe also calculated a second set of expected CLs assuming the predicted SM diboson yield in the PEs generation. In this case an excess in the observed CLs would indicate the presence of {\\em new physics}.\nCombining single--tagged and double--tagged results for all lepton categories, we find an expected limit of:\n\\begin{equation}\n\\left( 1.505 ^{+0.46}_{-0.61}\\right) \\times \\textrm{SM prediction.}\n\\end{equation}\n\nThe observed limit of $1.46$ times the SM prediction is thereby consistent with the existence of a signal and no presence of new physics.\n\nTable~\\ref{tab:limits} summarizes the expected (and observed) 95\\% production limits in units of the SM prediction. \n\n\\begin{table}\n\\begin{center}\n\\begin{tabular}{lccc}\n\\toprule\n Category & Expected 95\\% CL & Expected 95\\% CL & Observed \\\\ \n& (No Diboson) & (With Diboson) & 95\\% CL \\\\\\midrule\nSingle Tag &&&\\\\\nCEM+CMUP+CMX & $0.715 ^{+0.42}_{-0.38}$ & $1.625 ^{+0.54}_{-0.69}$ & 2.07 \\\\\nEMC & $1.215 ^{+0.57}_{-0.46}$ & $2.125 ^{+0.76}_{-0.6}$ & 1.70 \\\\ \\midrule\nDouble Tag &&&\\\\\nCEM+CMUP+CMX & $4.015 ^{+1.94}_{-1.50}$ & $4.875 ^{+2.15}_{-1.7}$ & 3.03 \\\\ \nEMC & $6.075 ^{+2.92}_{-2.20}$ & $7.015 ^{+3.15}_{-2.42}$ & 8.59 \\\\ \\midrule\nAll combined & $0.575 ^{+0.33}_{-0.31}$ & $1.505 ^{+0.46}_{-0.61}$ & 1.46 \\\\\n\\bottomrule \n\\end{tabular}\n\\caption[Expected and Observed 95\\% CL for Diboson Production]{Expected and observed 95\\% exclusion confidence levels for each lepton category, single and double tagged events, in units of the SM cross section for diboson production. Expected CLs are produced also including the diboson signal generated PEs, thus probing the contribution of new physics processes.}\\label{tab:limits}\n\\end{center}\n\\end{table}\n\n\n\n\\subsection{Significance and Cross Section Measurement}\n\nTo compute the significance of the signal, we performed a hypothesis test comparing the data to the likelihood ratio of the {\\em null} and {\\em test} hypotheses. \n\nThe {\\em null} hypothesis, $H_0$, assumes all the predicted processes except diboson production. The {\\em test} hypothesis, $H_1$, assumes that the diboson production cross section and the branching ratio into $HF$ are the ones predicted by the SM. The likelihood ratio is defined as:\n\\begin{equation}\\label{eq:l_ratio}\n-2 \\ln Q = -2 \\ln \\frac{ \\mathscr{L}(data|H_1,\\hat{\\theta})}{ \\mathscr{L}(data|H_0,\\hat{\\hat{\\theta}})}\n\\end{equation}\nwhere $\\mathscr{L}$ is defined by Equation~\\ref{eq:likelihood_full}, $\\theta$ represents the nuisance parameters describing the uncertain values of the quantities studied for systematic error, $\\hat{\\theta}$ the best fit values of $\\theta$ under $H_1$ and $\\hat{\\hat{\\theta}}$ are the best fit values of the nuisance parameters under $H_0$.\n\nTo perform the hypothesis test we generated two sets of PEs, one assuming $H_0$ and a second one assuming $H_1$ and we evaluated Equation~\\ref{eq:l_ratio} for each pseudo-data outcome. The distributions of the $- 2 \\ln Q$ values are shown in Figure~\\ref{fig:gaussians}. The integral of the $H_0$ $- 2 \\ln Q$ distribution below the $- 2 \\ln Q$ value of the real data gives the observed $p-$value of the signal.\n\nWe obtained an observed $p-$value of 0.00120, corresponding to a 3.03$\\sigma$ excess and producing the {\\em evidence} of the diboson signal. The result is compatible with the expectation as the {\\em test} hypothesis $p-$value, obtained from the median of the $H_1$ distribution of $ - 2 \\ln Q$, is $0.00126$.\n\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.95\\textwidth]{App7.5\/Figs\/combined_significance}\n\\caption[Significance Estimate of Diboson Signal with $ - 2 \\ln Q$]{Distributions of $ - 2 \\ln Q$ for the {\\em test} hypothesis $H_1$, which assumes the estimated backgrounds plus SM diboson production and decay into $HF$ (blue histogram), and for the {\\em null} hypothesis, $H_0$, which assumes no diboson (red histogram). The observed value of $ - 2 \\ln Q$ is indicated with a solid, vertical line and the $p$-value is the fraction of the $H_0$ curve integral to the left of the data.}\\label{fig:gaussians}\n\\end{figure}\n\nIn order to measure the diboson production cross section, a Bayesian marginalization technique is applied as described in Section~\\ref{sec:likelihood_meas}. The nuisance parameters are integrated, the maximum of the posterior distribution returns the cross section measurement while the smallest interval containing 68\\% of the integral gives the 1-$\\sigma$ confidence interval. The resulting cross section measurement is:\n\\begin{equation}\n \\sigma_{Diboson}^{Obs}=\\left(1.08^{+0.26}_{-0.40} \\right)\\times \\sigma_{Diboson}^{SM} = \\left(19.9^{+4.8}_{-7.4}\\right)\\textrm{~pb,}\n\\end{equation}\nwhere the errors include statistical and systematic uncertainties and $\\sigma^{SM}_{Diboson}$ is the SM predicted cross section derived from Table~\\ref{tab:mc_dib}\n\n\\chapter[SVM Multi-Jet Rejection]{Support Vector Machines Multi-Jet Rejection}\\label{chap:AppSvm}\n\nAn innovative multivariate method, based on the Support Vector Machines algorithm (SVM), is used in this thesis to drastically reduce the multi-jet background. \n\nOne of the crucial points in the search for diboson production in the $\\ell\\nu + HF$ final state is the maximization of the signal acceptance while keeping the background under control. This is a challenge because, in hadronic collider environment, jets are produced with a rate several order of magnitude larger than $W\\to \\ell \\nu$ events, therefore, as a jet can fake the lepton identification with not negligible probability (especially for electrons identification algorithms), multi-jet events are introduced in the sample.\n \nThe multi-jet background, a mixture of detector and physics processes, is challenging to parametrize and, usually, approximate data-driven models are obtained by appropriate fake-enriched selections (see Section~\\ref{sec:qcd}). These models are often statistically limited and the use of a different selection can produce unexpected biases in the simulated variables. It is obviously not trivial the use of multivariate techniques to tackle such a problem.\n\nThe SVM algorithm, described in Section~\\ref{sec:SVM}, is considered to perform well in this case as it offers good non-linear separation and stable solutions also on low statistical training samples~\\cite{svmbook, BShop_machine_learning}. It was never used before and we had to develop original solutions to address the major challenges: evaluate the robustness of the SVM against biases in the training set and establish the best, minimal set of input variables providing optimal performances. Section~\\ref{sec:method} describes how we solved the first problem while Section~\\ref{sec:var_sel} describes the input variable selection criteria.\n\nThe results, reported in Section~\\ref{sec:svm_res}, in terms of signal efficiency and background rejection, are superior to any other cut based or multi-variate method previously applied at CDF.\n\n\n\n\\section{Support Vector Machines}\\label{sec:SVM}\n\nThe SVM is a supervised learning binary classifier whose basic concept is the identification of the {\\em best separating hyper-plane} between two classes of $n$-dimension vectors. \n\nIn the case of linear separation the algorithm produces, given a training set of the vectors of the two classes, an {\\em unique} solution where the plane is defined by the minimum amount of vectors, called {\\em support vectors}, at the boundary of the two classes. In the case of non-linear separation, the plane is found in an abstract space, defined by a transformation of the input vectors. However it is not necessary to know the exact transformation, but just its effect on the scalar product between the vectors, named {\\em Kernel}, thus allowing a feasible solution. Finally the cases of not perfect separability of the two samples are included by introducing a penalty parameter accounting for the contamination. \n\nThe main advantages of the SVM with respect to other machine learning algorithms are the unique convergence of the problem, a small number of free tunable parameters (usually related to the Kernel choice) and good performances for low statistics training sets because only a small number of training vectors (the support vectors) are important for the final solution.\n\nIt is possible to find more details in~\\cite{svmbook, BShop_machine_learning}, but a short overview of the algorithm is also given in the following. For the actual, numerical, implementation of the SVM algorithm we relied on the \\texttt{LIBSVM} open source library~\\cite{libsvm}.\n\n\n\n\\subsection{The Linear Case}\nFigure~\\ref{fig:svm_linear} shows a basic example of the SVM linear classification separating two classes of bi-dimensional training vectors with a maximum margin hyperplane (a line for this simple case)\n\n\\begin{figure}\n \\begin{center} \n \\includegraphics[width=0.8\\textwidth]{AppSVM\/Figs\/vgv4}\n \\caption[Example of SVM Linear Separation]{An example of SVM: two linearly separable classes of data are represented with red and blue dots. The hyperplane (in this case a simple line) leading to a maximum margin separation is defined by the weight vector $w$ and the bias vector $b$.}\\label{fig:svm_linear}\n \\end{center}\n\\end{figure}\n\nThe problem can be formalized in a general way as the minimization of $|\\vec{w}|^2$ (with $\\vec{w}$ = vector normal to the plane) with the constraint:\n\\begin{equation}\\label{eq:constraint}\n y_i(\\vec{x_i}\\cdot\\vec{w} + b)-1\\geq 0\\qquad\n \\left\\{\\begin{array}{ll}\n y_i=+1; & i \\in \\textrm{signal} \\\\\n y_i=-1; & i \\in \\textrm{background}\\\\\n \\end{array} \\right.\n\\end{equation}\nThe problem has an unique solution obtained by the maximization of:\n\\begin{equation}\\label{eq:lagrange_1}\n L=\\sum_i \\alpha_i - \\frac{1}{2}\\sum_{i,j}\\alpha_i\\alpha_j y_i y_j\\vec{x_i}\\cdot\\vec{x_j} ,\n\\end{equation}\nobtained with the application of the Lagrange multipliers to Equation~\\ref{eq:constraint}.\nThe solution identifies, for some $i$, $\\alpha_i>0$. The associate vectors are the \\textit{support vectors}, i.e. a subset of the training sample that define the best hyper-plane (see Figure~\\ref{fig:svm_linear}).\n\nTo solve the case of not completely separable classes of vectors, a {\\em penalty parameter}, $C$, is added into the target function to account for the contamination. So that we have a new minimization condition:\n\\begin{equation}\n |\\vec{w}|^{2} + C \\sum_{i}{\\xi_{i}};\n\\end {equation}\nand a new constraint (derived again from Equation~\\ref{eq:constraint}):\n\\begin{equation}\\label{eq:lagrange_2}\n y_i(\\vec{x_i}\\cdot\\vec{w} + b) \\geq 1-\\xi_{i}\\quad\\textrm{with}\\quad \\xi \\geq 0\\textrm{.}\n\\end{equation}\n\nThe parameter $C$ defines, before the training, the SVM implementation therefore it represents as one of the {\\em hyper parameter} of the SVM.\n\nA newly seen vector, $\\vec{X}$, is classified according the position with respect to the plane defined by the support vectors $\\vec{x_i}$ and the parameters $\\alpha_i$:\n\\begin{equation}\\label{eq:distance}\n D(\\vec{X})= \\sum_i\\alpha_i y_i \\vec{x_i}\\cdot \\vec{X_i} - b,\n\\end{equation}\nwhere $b$ is a bias term of the solution. The sign of $D(\\vec{X})$ defines the classification but the value itself can be seen as the distance of a test vector $\\vec{X}$ from the classification plane. However, as we will see in the next paragraph, a non-linear classification is possible only thanks to a not-explicit transformation in a different vector space, where $D$ looses it immediate geometrical meaning. \n\nCommonly SVMs are used as binary classifiers but, here, we add a large degree of flexibility by exploiting the full information of the variable $D$. We see it as a dimensionality reducer that summarizes all the information obtained during the training and classification process. \n\n\n\n\\subsection{Kernel Methods}\n\nNon-linearly separable classes of vectors can be classified by transforming them into linearly separable\nclasses. An opportune function, $\\Phi(\\vec{x})$, can be used to map the elements into another\nspace, usually with higher dimension where the separation is possible. \n\nHowever the identification of $\\Phi(\\vec{x})$ is {\\em non trivial} and the, so called, {\\em Kernel trick} is often used: a Kernel function, $\\mathbf{K}(x_i,x_j)$, generalizes the scalar product appearing in Equation~\\ref{eq:lagrange_1} (or Equation~\\ref{eq:lagrange_2}) without the need of explicitly know $\\Phi(\\vec{x})$. Or in equations, we compose the mapping $\\Phi(\\vec{x})$ with the inner product:\n\\begin{equation}\n \\mathbf{K}(x_i,x_j)= \\Phi(x_i)\\cdot\\Phi(x_j)\\quad\\mathrm{with}\\quad \\Phi : \\Re^n\\mapsto \\mathcal{H}.\n\\end{equation}\n\nThe function $\\mathbf{K}$ should satisfy to a general set of roles to be a Kernel, but we want only to briefly describe the {\\em Gaussian} Kernel we used in this work. It is expressed as:\n\\begin{equation}\\label{eq:gauss_k}\nK(x_{i},x_{j}) = e^{-\\gamma |\\vec{x}_{i}-\\vec{x}_{j}|^2}\n\\end{equation}\nThe corresponding $\\Phi(x)$ maps to an infinite dimension space and it is not known. The Kernel is defined only by one hyper-parameter, $\\gamma$, that should be defined before the training.\n\n\n\n\n\n\\section{SVM Training in a Partially Biased Sample}\\label{sec:method}\n\nThe assumption behind the supervised learning is that the labelled samples, used for the classifier training, are drawn from the same probability distribution of the unclassified events. \nHowever in our case of study, where only an approximate and statistically limited model of the background processes is available (see Sections~\\ref{sec:cdf_set} and~\\ref{sec:qcd} for the multi-jet background description), we do not expected the previous assumption to hold for every region of the phase space. To cope with this problem, we developed an original methodology to evaluate the SVM training performances.\n\nSection~\\ref{sec:SVM} shows that, for each choice of hyper-parameters and training vectors, only one optimal SVM solution exists and, for it, we need to evaluate the performances. \n\nAs a performances estimator we use the {\\em confusion matrix} of the classifier: the element $(i,j)$ of the matrix is the fraction of the class $i$ classified as member of class $j$. Figure~\\ref{fig:confusion} shows a representation of it in the two classes case, where one class is considered the {\\em background} and the other the {\\em signal}. We obtain a reliable estimate of the classifier quality by filling the confusion matrix in two independent ways and combining all the available information.\n\n\\begin{figure}[h!]\n \\begin{center} \n \\begin{tabular}{|c|c|}\n \\toprule\n {\\em Sgn} classified as {\\em Sgn} & {\\em Bkg} classified as {\\em Sgn}\\\\\n \\midrule\n {\\em Sgn} classified as {\\em Bkg} & {\\em Bkg} classified as {\\em Bkg}\\\\\n \\bottomrule\n \\end{tabular} \n \\caption[Definition of Confusion Matrix]{Definition of confusion matrix for a two classes ($Sgn$ and $Bkg$) classification problem. This reproduces the case of an algorithm used to discriminate signal {\\em vs} background: the elements of the matrix are the signal and background classification performances and the cross contamination.}\\label{fig:confusion}\n \\end{center}\n\\end{figure}\n\nThe first performance evaluation method is the {\\em $k$-fold cross-validation}: the training set is divided into $k$ sub-samples of which one is used as a validation set and the remaining $k - 1$ are used in the training; the confusion matrix is then evaluated applying the trained discriminant to the validation set. The cross-validation process is repeated $k$ times, the {\\em folds}, and the final performance is given by the average on all the folds. This method is solid against over fitting but it has no protection against biases on the complete training sample.\n\nThe second method, a {\\em key feature} of this work, is based on a bi-component fit that uses signal and background templates and it is performed on a significant distribution of the unclassified events, the {\\em data}. While the signal and background templates are derived in the same way of the training samples, the unclassified data events are, by definition composed by an unknown mixture of the {\\em true} signal and background events. The fit is performed by maximizing a binned likelihood function, $\\lambda$, where the Poisson statistic of the templates is used and the fractions of the signal and of background templates, from which we can derive the elements of the confusion matrix, are free parameters. The fitting function is implemented in the ROOT~\\cite{rootr} analysis package and it is derived from~\\cite{Barlow_1993}. \nFigure~\\ref{fig:toysamplefit} shows an example of the fit used on the toy model described in the next Section.\n\nIf the variable considered in the fit is not well reproduced in the simulation then we expect that the fitted fractions will differ greatly from the results obtained with the $k$-fold cross-validation. At the same time we can evaluate quantitatively the agreement between the data shape and the fitted templates because the quantity:\n\\begin{equation}\n \\chi^2 = -2\\ln(\\lambda),\n\\end{equation}\nfollows a $\\chi^2$ probability distribution (under general assumptions).\n\nThe last critical point is the identification of a {\\em sensitive} variable to be used in the fit. In a previous work~\\cite{CHEP_svm} we exploited the \\mbox{${\\not}{E_T}$} distribution as it is sensitive to the multi-jet contamination. A much more general approach, by the machine learning point of view, is the use of the SVM distance value, $D$, defined in Equation~\\ref{eq:distance}. If the SVM training performances are optimal, also the variable $d$ offers an optimal discrimination, furthermore the cross check on the $\\chi^2$ of the template fit ensure a good shape agreement between the data and signal and background templates. We verified the validity of the fit procedure with a toy example reported in the following.\n\n\\begin{figure}[!h]\n \\begin{center}\n \\includegraphics[width=0.9\\textwidth]{AppSVM\/Figs\/toysample}\n \\caption[Bi-component Template Fit on Toy Data]{A bi-component fit is performed on the SVM distance, $D$ (Equation~\\ref{eq:distance}), of a signal (blue) and background (red) templates toy data.}\\label{fig:toysamplefit} \n \\end{center}\n\\end{figure}\n \n\n\\subsection{A Toy Example}\n\nWe built a toy example in order to verify the robustness of the proposed method for an SVM performances evaluation when partially biased samples are available.\n\nThe toy is composed by three data-sets generated with known probability distributions:\n\\begin{description}\n\\item[signal model:] $10^5$ vectors generated from a $2-$Dim Gaussian distribution with the following mean, $\\vec{\\mu}_{Sgn}$ and standard deviation, $\\tilde{\\sigma}_{Sgn}$:\n \\begin{equation}\\label{eq:sgn_gaus}\n \\vec{\\mu}_{Sgn} = \\left[\n \\begin{array}{c}\n - 3 \\\\\n 0 \\end{array} \\right]\n ,\\qquad \n \\tilde{\\sigma}_{Sgn} = \\left[\n \\begin{array}{cc}\n 8 & 0 \\\\\n 0 & 8 \\end{array} \\right].\n \\end{equation}\n \n \\item[Background model:] $10^5$ vectors generated from a $2-$Dim Gaussian distribution with the following mean, $\\vec{\\mu}_{Bkg}$ and standard deviation, $\\tilde{\\sigma}_{Bkg}$:\n \n \\begin{equation}\\label{eq:bkg_gaus}\n \\vec{\\mu}_{Bkg} = \\left[ \n \\begin{array}{c}\n 3 \\\\\n 0 \\end{array} \\right]\n ,\\qquad \n \\tilde{\\sigma}_{Bkg} = \\left[\n \\begin{array}{cc}\n 8 & 0 \\\\\n 0 & 8 \\end{array} \\right].\n \\end{equation}\n \\item[Data:] a mixture of $5 \\cdot 10^4$ vectors generated from the same distribution of the signal model (Equation~\\ref{eq:sgn_gaus}) and $5\\cdot 10^4$ vectors generated from a {\\em true} background distribution similar to the background model (Equation~\\ref{eq:bkg_gaus}) but with $\\tilde{\\sigma}_{Bkg}$ increased by 20\\% in one direction to simulate a mismatch between the real background and the model.\n \\end{description}\n %\n We tested several combinations of the hyper-parameters $C$ and $\\gamma$ (over a grid) using the signal and background model in the training. For obtained SVM we evaluated the $k$-fold cross validation and we performed the template fit on the SVM distance, $D$, evaluated on the data sample. The result is reported in Figure~\\ref{fig:ToyResult} with the real performances reported on the $x$ axis of the diagram (we know the true label of the data vectors). The evaluation of the performances obtained with the fit is on the diagonal of the plane, therefore it gives a much more realistic estimate of the true performances of the classifier.\n\n \\begin{figure}[!h]\n \\begin{center}\n \\includegraphics[width=0.75\\textwidth,height=0.5\\textwidth]{AppSVM\/Figs\/ToyResult}\n \\caption[Toy Model SVM Performances Estimate with Fit and Cross-Validation]{SVM performances estimate with a $k$-fold cross-validation (green circles) and with a bi-component signal and background template fit on the SVM distance (red triangle) of toy data of known composition. The {\\em true} performances of the SVM classifier are reported on the $x$ axis. The fit evaluation appears on the diagonal of the plane, signaling a more realistic estimate of the true performances of the classifier. }\\label{fig:ToyResult} \n \\end{center}\n \\end{figure}\n\n\n\n\\section{Performances on the CDF Dataset}\\label{sec:cdf_set}\n\nThe final goal of the SVM discriminant we discussed is the realization of a tool able to reject the multi-jet background in a wide range of searches performed in the lepton plus jets channel in a hadron collider environment. \n\nWe performed the SVM training on $W\\to e\\nu + $jets candidate, as the electron identification is more tamed by the multi-jet background. Furthermore we performed the training process two times, one in the central ($|\\eta|<1.1$) and one in the forward ($1.2<|\\eta|<2.0$) region of the detector, as both the electron identification algorithm (see Section~\\ref{sec:lep_id}) and the kinematic of the events are different. \n\n\n\\subsection{Training Sample Description}\\label{sec:training}\n\nWe defined both a central and a forward training set using $7000$ $W\\to e\\nu +$jets signal events and $3500$ multi-jet background events:\n\\begin{description}\n \\item[Signal:] $W + 2,3$ partons ALPGEN~\\cite{alpgen} MC, where the $W$ is forced to decay into electron and neutrino. We have about $10^5$ generated events and we keep approximately $9\\times 10^4$ events as a control sample (i.e. not used for training). The CEM and PHX electron identification algorithms are used for the central and forward sample selection. \n\\item[Background:] we obtain a suitable background sample with a data-driven approach. The anti-electrons selection described in Section~\\ref{sec:qcd} is used for the central background training while for the forward training we had study a improved multi-jet model. \n\nIn particular, we noticed that the anti-PHX sample produced large over estimates in the detector region of $1.2<|\\eta|<1.4$ where two different calorimeter sub-systems are connected. These events are clear fake leptons and they are usually rejected by very loose kinematic requirements, nevertheless the quantity of them produced a too large bias in the SVM training efficiency estimate. Non-Isolated PHX electrons (with $IsoRel>0.1$) were found to give a reliable training set after a correction to the lepton $E_T$ equal to the amount of energy in the outer isolation cone ($E_T^{0.1<\\Delta R<0.4}$).\n\\item[Data:] The data sample used for the bi-component fit validation corresponds to the data periods (see Table~\\ref{tab:datasets}) from p18 to p24. They represent only a small fraction of the total CDF data with intermediate luminosity profile and run conditions.\n\\end{description}\nA close to final lepton plus jet selection with loose kinematic cuts is applied to all the samples. We require two or more jets reconstructed (see Section~\\ref{sec:jetCone}) in the central region of the detector ($\\eta<2.0$), with energy corrected at L5 and $E_T> 18$~GeV, also a minimal amount of \\mbox{${\\not}{E_T}$} $>15$~GeV is used as signature of the escaping neutrino. We also removed events with not understood behaviour like electron $E_T>300$~GeV, $M_T^W>200$~~GeV$\/c^{2}$ and, only for the forward sample, \\mbox{${\\not}{E_T^{raw}}$}$<20$~GeV (this last cut was needed to avoid trigger turn-on effects).\n\n\n\n\\subsection{Variable Selection}\\label{sec:var_sel}\n\nA multivariate algorithm relies on a given set of input variables. The {\\em feature selection} problem is fundamental in machine-learning and, if possible, even more in the present case where the background sample does not guaranteed a perfectly model of all the variables. \n\nWe started from a large set of twenty-four variables chosen according to two basic criteria: no correlation with respect to the lepton identification variables and exploit of the kinematic difference between real $W+$ jets events and multi-jet fakes. These requirements allowed the development of a flexible multi-jet rejection algorithm, applied with very good performances also to muon and isolated track lepton selections.\n\nTable~\\ref{tab:input_vars} shows all the input variables that we used in the optimization process. Many of them were introduced in Chapter~\\ref{chap:objects} but the following are new:\n\\begin{itemize}\n\\item \\mbox{${\\not}{p_T}$} is the missing momentum defined as the momentum imbalance on the transverse plane. It is computed adding all the reconstructed charged tracks transverse momenta, $\\vec{p_i}$:\n \\begin{equation}\n \\vec{{\\not}{p_T}}\\equiv - \\sum_i{\\vec{p^i}_T}\\quad\\textrm{with }|\\vec{p^i}_T|>0.5\\textrm{GeV}\/c\\textrm{;}\n \\end{equation}\n\\item $M_T^W$ is the {\\em transverse mass} of the reconstructed $W$ boson:\n \\begin{equation}\n M_T^W=\\sqrt{2(E_T^{lep}\\cancel{E}_T - E_x^{lep}\\cancel{E}_x - E_y^{lep}\\cancel{E}_y)}\\mathrm{.}\n \\end{equation}\n\\item $MetSig$ is the \\mbox{${\\not}{E_T}$} {\\em significance}, a variable that relates the reconstructed \\mbox{${\\not}{E_T}$} with the detector activity (jets and unclustered energy):\n \\begin{equation}\n MetSig =\\frac{\\cancel{E}_T}{\\sqrt{ \\Delta E^{jets}+ \\Delta E^{uncl}}}\\mathrm{,}\n \\end{equation}\n where:\n \\begin{eqnarray}\n \\Delta E^{jets} =\\sum_{j}^{jets}(cor_{j}^2\\cos^2\\Big(\\Delta\\phi\\big(\\vec{p_j},\\cancel{\\vec{E}}_T\\big)\\Big)E_{T}^{raw, j}\\mathrm{,}\\\\\n \\Delta E^{uncl} = \\cos^2\\Big(\\Delta\\phi\\big(\\vec{E}_T^{uncl}\\cancel{\\vec{E}}_T\\big)\\Big)E_T^{uncl}\\mathrm{,}\n \\end{eqnarray}\n $uncl$ refers to the calorimeter energy not clustered into electrons or jets and $cor_j$ is the total correction applied to each jet.\n\\item $\\nu^{Min},\\quad \\nu^{Max}$ are the two possible reconstruction of the neutrino momenta. As the $p_z^{\\nu}$ component is not directly measurable we infer it from the $W$ boson mass and the lepton momentum. The constraints lead to a quadratic equation which may have two real solutions, one real solution, or two complex solutions\\footnote{The real part is chosen in this case}. The reconstructed $\\nu^{Min}$, $\\nu^{Max}$ derive from the distinction of $p_z^{\\nu,Max}$ and $p_z^{\\nu,Min}$.\n\\end{itemize}\n\n\\begin{table}\n \\begin{tabular}{cl| cl| cl| cl }\n \\toprule\n \\multicolumn{8}{c}{Possible Input Variables }\\\\\n \\midrule\n 1 & $p_T^{lep}$ & 7 & $E_T^{raw, jet1} $ & 13 & $\\Delta\\phi( {\\not}{p_T},$ $lep)$ & 19 & $\\Delta R( lep,$ $jet 2)$ \\\\\n 2 & \\mbox{${\\not}{E_T}$} & 8 & $E_T^{raw, jet2} $ & 14 & $\\Delta\\phi( {\\not}{p_T},$ \\mbox{${\\not}{E_T}$} $)$ & 20 & $\\Delta R( \\nu^{min},$ $jet 1)$ \\\\ \n 3 & \\mbox{${\\not}{E_T^{raw}}$} & 9 & $E_T^{cor,jet1} $ & 15 & $\\Delta\\phi( {\\not}{p_T},$ \\mbox{${\\not}{E_T^{raw}}$}$)$ & 21 & $\\Delta R( \\nu^{min},$ $jet 2)$ \\\\ \n 4 & \\mbox{${\\not}{p_T}$} & 10 & $E_T^{cor,jet2} $ & 16 & $\\Delta\\phi( lep,$ \\mbox{${\\not}{E_T}$} $)$ & 22 & $\\Delta R( \\nu^{min},$ $lep)$ \\\\ \n 5 & $M_T^W$ & 11 & $\\Delta\\phi( jet 1,$ \\mbox{${\\not}{E_T}$} $)$ & 17 & $\\Delta\\phi( lep,$ \\mbox{${\\not}{E_T^{raw}}$}$)$ & 23 & $\\Delta R( \\nu^{max},$ $jet 1)$ \\\\ \n 6 & $MetSig$ & 12 & $\\Delta\\phi( jet 2,$ \\mbox{${\\not}{E_T}$} $)$ & 18 & $\\Delta R( lep,$ $jet 1)$ & 24 & $\\Delta R( \\nu^{max},$ $jet 1)$ \\\\ \n\\bottomrule\n \\end{tabular}\n \\caption[Possible SVM Input Variables]{All the possible input variables used for the SVM training and optimization. See Section~\\ref{sec:var_sel} for a detailed description.}\\label{tab:input_vars} \n\\end{table}\n\nUnluckily the extensive research over all the possible combinations of variables across all the $C, \\gamma$ phase space of a given SVM training, is computationally unfeasible. To scan the most relevant sectors of the phase space we applied factorized and incremental optimization:\n\\begin{itemize}\n\\item for all the configuration of {\\em three} variables and the given training set, we evaluate a grid of $C, \\gamma$ values in the intervals\\footnote{The use of a logarithmic scale allows to scan the parameters across different orders of magnitude.}:\n \\begin{equation}\n \\log_2 C\\in [-3,8] \\quad\\textrm{and}\\quad \\log_2\\gamma\\in[-4,5].\n \\end{equation}\nWe select only the best training configuration according to the confusion matrix evaluation.\n\\item For each {\\em best} SVM of a given variable configuration we perform a bi-component fit on the SVM distance $D$. We evaluate the $\\chi^2$ of the fit, reduced by the Number of Degrees of Freedom ($NDoF$), and we compare the fitted background contamination, $f_{Bkg}^{Fit}$, against the one obtained from the $n$-fold cross-validation, $f_{Bkg}^{n-fold}$. The SVM under exam is rejected if:\n \\begin{equation}\n \\frac{\\chi^2}{NDoF}>3\\quad or \\quad\\frac{f_{Bkg}^{Fit}}{f_{Bkg}^{n-fold}}>2\n \\end{equation}\n\n\\item The remaining SVMs are displayed on a {\\em signal-efficiency vs background-contamination} scatter plot like the one in Figure~\\ref{fig:discriminants}. The 5 best variable combinations are \nselected for further processing.\n\\item We add other 2 or 3 variables to the best variables combinations obtained in the previous step and we iterate the chain.\n\\end{itemize}\nAfter a couple of iterations the best variable combination and $C,\\gamma$ hyper-parameters choice remains stable within $1\\div 2\\%$.\n\n\\begin{figure}[!h]\n \\begin{center} \n \\includegraphics[width=0.9\\textwidth]{AppSVM\/Figs\/discriminants_bests_trainCEM3}\n \\caption[Signal Efficiency {\\em vs} Background Contamination for Multiple SVM]{Different SVM configurations (i.e. with different input variables), obtained for the central region training, are displayed on a {\\em signal-efficiency vs background-contamination} scatter plot. The signal efficiency is directly estimated from simulation while the background contamination is obtained from the bi-component template fit of the SVM distance, $D$, described in Section~\\ref{sec:method}. The blue starts represent the performances obtained from a three (out of twenty-four) input variables training. The best five configurations w.r.t. the Euclidean distance from the optimal point ($\\epsilon_{Sig}=1$, $f_{Bkg}=0$) are circled in yellow. The Euclidean distance of the fifth best SVM is represented by a dotted line. Three iterations, with {\\em three, six and nine} input variables are represented by dotted lines, after that no more appreciable improvement occurs.}\\label{fig:discriminants}\n \\end{center}\n\\end{figure}\n\n\n\\section{Final SVM Results}\\label{sec:svm_res}\nThe SVM configurations obtained by the process described in this appendix are finally used in the multi-jet rejection phase of the analysis (Section~\\ref{sec:svm_sel}) and in the background normalization estimates\\footnote{By construction a fit on the SVM distance, $D$, offers a reliable estimate of the multi-jet background normalization.} (see Sections~\\ref{sec:wjets_pretag} and~\\ref{sec:qcd}).\n\nThe optimal hyper-parameter configurations that we obtain for the central (superscript $c$) and forward (superscript $f$) SVMs are:\n\\begin{equation}\nC^c = 7,\\quad \\gamma^c = \u2212- 1;\n\\end{equation}\n\\begin{equation}\nC^f = 8,\\quad \\gamma^f = \u2212- 1.\n\\end{equation}\nTable~\\ref{tab:svm_fin_vars} reports the final eight input variables used for the central SVM and the six ones used for the forward SVM.\n\\begin{table}\n \\begin{center} \n \\begin{tabular}{lccc}\n \\toprule\n \\multicolumn{4}{c}{ Final SVM Input Variables}\\\\\n \\midrule\n Central SVM: & $M^W_T$ & \\mbox{${\\not}{E_T^{raw}}$} & ${\\not}{p_T}$\\\\ \n & $MetSig$ & $\\Delta \\phi({\\not}{p_T},$ \\mbox{${\\not}{E_T}$} ) & $\\Delta \\phi(lep,$ \\mbox{${\\not}{E_T}$} $)$ \\\\\n & $\\Delta R(\\nu^{Min},$ $lep)$ & $\\Delta \\phi(Jet1,$ \\mbox{${\\not}{E_T}$} ) & \\\\\n \\midrule\n Forward SVM: & $M^W_T$ & \\mbox{${\\not}{E_T^{raw}}$} & ${\\not}{p_T}$ \\\\\n & $MetSig$ & $\\Delta \\phi({\\not}{p_T},$ \\mbox{${\\not}{E_T}$} ) & $\\Delta \\phi({\\not}{p_T},$ \\mbox{${\\not}{E_T^{raw}}$}) \\\\\n \\bottomrule\n \\end{tabular}\n \\caption[Final SVM Input Variables]{Input variables used for the configuration of the {\\em central} and {\\em forward} SVM multi-jet discriminants.}\\label{tab:svm_fin_vars}\n \\end{center}\n\\end{table}\nFigure~\\ref{fig:final_svm} shows the complete shape of the two discriminants for the multi-jet background models and the $W+2$ partons signal. A quantitative measurement of the performances can be seen in Table~\\ref{tab:svm_eff}.\n\n\\begin{figure}[!h]\n \\begin{center} \n \\includegraphics[width=0.495\\textwidth]{AppSVM\/Figs\/central_svm}\n \\includegraphics[width=0.495\\textwidth]{AppSVM\/Figs\/forward_svm}\n \\caption[Final SVM Discriminants for Central and Forward Detector Regions]{Distribution of the SVM distance, $D$, described in Equation~\\ref{eq:distance} for the central (left) and the forward (right) SVM discriminants obtained from the optimization process. Multi-jet background models are shown in red, $W+2$ partons MC signal is shown in blue.}\\label{fig:final_svm}\n \\end{center}\n\\end{figure}\n\nWe can conclude that we successfully built a multi-jet rejection tool based of the SVM algorithm, nevertheless the challenges of a multi-variate approach to this problem. The CDF II dataset was a perfect test-bench for this problem with very good performances, however the procedure can be exported also to any of the LHC experiments.\n\n\n\\chapter{Rate Systematics Summary}\\label{app:rate_sys}\n\nTables from~\\ref{tab:1tag_cem} to~\\ref{tab:2tag_emc}: rate systematics variations for each channel entering in the statistical analysis (see Chapter~\\ref{chap:StatRes}), same name systematics are fully correlated.\n\n\\renewcommand{\\arraystretch}{0.97}\n\\renewcommand{\\tabcolsep}{1.5pt}\n\n\\begin{table}[!h]\n \\begin{scriptsize}\n \\begin{tabular}{ c c c c c c c c c c c c } \\toprule\n & nonW & $t\\bar{t}$ & s-top $s$ & s-top $t$ & Zjets & W+LF & Wbb & Wcc & Wcj & WW & WZ\/ZZ \\\\\\midrule\n$\\texttt{ XS\\_ttbar}$ & $0\/0$ & $10\/-10$ & $10\/-10$ & $10\/-10$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFBTAGSF}$ & $0\/0$ & $3\/-3$ & $3\/-3$ & $5\/-5$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $11\/-11$ & $5\/-5$ \\\\\n$\\texttt{ CDFLUMI}$ & $0\/0$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $6\/-6$ & $6\/-6$ \\\\\n$\\texttt{ LEPACC\\_CEM}$ & $0\/0$ & $2\/-2$ & $2\/-2$ & $2\/-2$ & $2\/-2$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $2\/-2$ & $2\/-2$ \\\\\n$\\texttt{ QCD\\_CEM}$ & $40\/-40$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFVHF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $30\/-30$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFWC}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFMISTAG}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $11\/-11$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ZJETS}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $45\/-45$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ISRFSRPDF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $4\/-4$ & $4\/-4$ \\\\\n$\\texttt{ CDFJES}$ & $0\/0$ & $-12\/12$ & $-3\/2$ & $-2\/1$ & $1\/6$ & $1\/-1$ & $-6\/6$ & $-3\/6$ & $-5\/6$ & $-3\/-4$ & $2\/-2$ \\\\\n$\\texttt{ CDFQ2}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $-2\/-1$ & $4\/11$ & $5\/8$ & $3\/9$ & $0\/0$ & $0\/0$ \\\\\n\\hline\n$\\texttt{TOT}$ & $40\/40$ & $17\/17$ & $12\/12$ & $13\/13$ & $45\/46$ & $11\/11$ & $31\/33$ & $31\/32$ & $30\/32$ & $14\/14$ & $9\/11$ \\\\\n\\hline\n \\end{tabular}\n \\end{scriptsize}\n \\caption[Rate Uncertainties: Single-Tag, CEM Channel]{Per-sample rate uncertainty (\\% up\/down). Single-tag, CEM channel.}\\label{tab:1tag_cem}\n \\end{table}\n\n \\begin{table}[!h]\n \\begin{scriptsize}\n \\begin{tabular}{ c c c c c c c c c c c c} \\toprule\n & nonW & $t\\bar{t}$ & s-top $s$ & s-top $t$ & Zjets & W+LF & Wbb & Wcc & Wcj & WW & WZ\/ZZ \\\\\\midrule\n$\\texttt{ XS\\_ttbar}$ & $0\/0$ & $10\/-10$ & $10\/-10$ & $10\/-10$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFBTAGSF}$ & $0\/0$ & $3\/-3$ & $3\/-3$ & $5\/-5$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $11\/-11$ & $5\/-5$ \\\\\n$\\texttt{ CDFLUMI}$ & $0\/0$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $6\/-6$ & $6\/-6$ \\\\\n$\\texttt{ LEPACC\\_PHX}$ & $0\/0$ & $2\/-2$ & $2\/-2$ & $2\/-2$ & $2\/-2$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $2\/-2$ & $2\/-2$ \\\\\n$\\texttt{ QCD\\_PHX}$ & $40\/-40$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFVHF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $30\/-30$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFWC}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFMISTAG}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $11\/-11$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ZJETS}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $45\/-45$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ISRFSRPDF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $4\/-4$ & $4\/-4$ \\\\\n$\\texttt{ CDFJES}$ & $0\/0$ & $-8\/12$ & $-1\/2$ & $-3\/3$ & $2\/-6$ & $-1\/-1$ & $-4\/4$ & $-5\/4$ & $-8\/6$ & $-7\/-5$ & $-10\/-3$ \\\\\n$\\texttt{ CDFQ2}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $-1\/-2$ & $0\/15$ & $0\/18$ & $2\/10$ & $0\/0$ & $0\/0$ \\\\\n\\hline\n$\\texttt{TOT}$ & $40\/40$ & $14\/17$ & $12\/12$ & $13\/13$ & $45\/46$ & $11\/11$ & $30\/34$ & $30\/35$ & $31\/32$ & $15\/14$ & $14\/10$ \\\\\n\\hline\n\\end{tabular}\n \\end{scriptsize}\n \\caption[Rate Uncertainties: Single-Tag, PHX Channel]{Per-sample rate uncertainty (\\% up\/down). Single-tag, PHX channel.}\\label{tab:1tag_phx}\n \\end{table}\n\n \\begin{table}\n \\begin{scriptsize}\n\n\\begin{tabular}{ c c c c c c c c c c c c } \\toprule\n & nonW & $t\\bar{t}$ & s-top $s$ & s-top $t$ & Zjets & W+LF & Wbb & Wcc & Wcj & WW & WZ\/ZZ \\\\\\midrule\n$\\texttt{ XS\\_ttbar}$ & $0\/0$ & $10\/-10$ & $10\/-10$ & $10\/-10$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFBTAGSF}$ & $0\/0$ & $3\/-3$ & $3\/-3$ & $5\/-5$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $11\/-11$ & $5\/-5$ \\\\\n$\\texttt{ CDFLUMI}$ & $0\/0$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $6\/-6$ & $6\/-6$ \\\\\n$\\texttt{ LEPACC\\_MU}$ & $0\/0$ & $1\/-1$ & $1\/-1$ & $1\/-1$ & $1\/-1$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $1\/-1$ & $1\/-1$ \\\\\n$\\texttt{ QCD\\_MU}$ & $40\/-40$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFVHF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $30\/-30$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFWC}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFMISTAG}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $11\/-11$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ZJETS}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $45\/-45$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ISRFSRPDF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $4\/-4$ & $4\/-4$ \\\\\n$\\texttt{ CDFJES}$ & $0\/0$ & $-11\/12$ & $-4\/3$ & $-1\/1$ & $9\/-6$ & $1\/1$ & $-8\/6$ & $-5\/5$ & $-6\/7$ & $-1\/-4$ & $1\/-2$ \\\\\n$\\texttt{ CDFQ2}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $-3\/-4$ & $2\/16$ & $6\/12$ & $2\/15$ & $0\/0$ & $0\/0$ \\\\\n\\hline\n$\\texttt{TOT}$ & $40\/40$ & $17\/17$ & $13\/12$ & $13\/13$ & $46\/46$ & $11\/12$ & $31\/35$ & $31\/33$ & $31\/34$ & $13\/14$ & $9\/16$ \\\\\n\\hline\n\\end{tabular}\n \\end{scriptsize}\n\\caption[Rate Uncertainties: Single-Tag, CMUP $+$ CMX Channel]{Per-sample rate uncertainty (\\% up\/down). Single-tag, CMUP $+$ CMX channel.}\\label{tab:1tag_emc}\n \\end{table}\n\n\n \\begin{table}\n \\begin{scriptsize}\n\\begin{tabular}{ c c c c c c c c c c c c } \\toprule\n & nonW & $t\\bar{t}$ & s-top $s$ & s-top $t$ & Zjets & W+LF & Wbb & Wcc & Wcj & WW & WZ\/ZZ \\\\\\midrule\n$\\texttt{ XS\\_ttbar}$ & $0\/0$ & $10\/-10$ & $10\/-10$ & $10\/-10$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFBTAGSF}$ & $0\/0$ & $3\/-3$ & $3\/-3$ & $5\/-5$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $11\/-11$ & $5\/-5$ \\\\\n$\\texttt{ CDFLUMI}$ & $0\/0$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $6\/-6$ & $6\/-6$ \\\\\n$\\texttt{ LEPACC\\_EMC}$ & $0\/0$ & $5\/-5$ & $5\/-5$ & $5\/-5$ & $5\/-5$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $5\/-5$ & $5\/-5$ \\\\\n$\\texttt{ QCD\\_EMC}$ & $40\/-40$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFVHF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $30\/-30$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFWC}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFMISTAG}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $11\/-11$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ZJETS}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $45\/-45$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ISRFSRPDF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $4\/-4$ & $4\/-4$ \\\\\n$\\texttt{ CDFJES}$ & $0\/0$ & $-10\/9$ & $0\/-2$ & $0\/-3$ & $11\/-8$ & $0\/-0$ & $-5\/6$ & $-8\/4$ & $-3\/2$ & $9\/-6$ & $8\/-6$ \\\\\n$\\texttt{ CDFQ2}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $-3\/-2$ & $11\/13$ & $12\/11$ & $11\/16$ & $0\/0$ & $0\/0$ \\\\\n\\hline\n$\\texttt{TOT}$ & $40\/40$ & $17\/16$ & $13\/13$ & $14\/14$ & $47\/46$ & $11\/11$ & $32\/33$ & $33\/32$ & $32\/34$ & $17\/15$ & $14\/19$ \\\\\n\\hline\n\\end{tabular}\n \\end{scriptsize}\n\\caption[Rate Uncertainties: Single-Tag, EMC Channel]{Per-sample rate uncertainty (\\% up\/down). Single-tag, EMC channel.}\\label{tab:1tag_mu}\n \\end{table}\n\n\n \\begin{table}\n \\begin{scriptsize}\n\\begin{tabular}{ c c c c c c c c c c c c } \\toprule\n & nonW & $t\\bar{t}$ & s-top $s$ & s-top $t$ & Zjets & W+LF & Wbb & Wcc & Wcj & WW & WZ\/ZZ \\\\\\midrule\n$\\texttt{ XS\\_ttbar}$ & $0\/0$ & $10\/-10$ & $10\/-10$ & $10\/-10$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFBTAGSF}$ & $0\/0$ & $11\/-11$ & $11\/-11$ & $12\/-12$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $24\/-24$ & $11\/-11$ \\\\\n$\\texttt{ CDFLUMI}$ & $0\/0$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $6\/-6$ & $6\/-6$ \\\\\n$\\texttt{ LEPACC\\_CEM}$ & $0\/0$ & $2\/-2$ & $2\/-2$ & $2\/-2$ & $2\/-2$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $2\/-2$ & $2\/-2$ \\\\\n$\\texttt{ QCD\\_CEM}$ & $40\/-40$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFVHF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $30\/-30$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFWC}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFMISTAG}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $21\/-21$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ZJETS}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $45\/-45$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ISRFSRPDF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $4\/-4$ & $4\/-4$ \\\\\n$\\texttt{ CDFJES}$ & $0\/0$ & $-4\/10$ & $-5\/3$ & $-4\/2$ & $-9\/-1$ & $3\/-3$ & $-7\/8$ & $-15\/2$ & $-11\/9$ & $-18\/6$ & $-4\/-5$ \\\\\n$\\texttt{ CDFQ2}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/1$ & $18\/9$ & $6\/17$ & $5\/-9$ & $0\/0$ & $0\/0$ \\\\\n\\hline\n$\\texttt{TOT}$ & $40\/40$ & $17\/19$ & $17\/16$ & $17\/17$ & $46\/45$ & $21\/21$ & $35\/32$ & $34\/35$ & $32\/32$ & $31\/26$ & $14\/14$ \\\\\n\\hline\n\\end{tabular}\n \\end{scriptsize}\n\\caption[Rate Uncertainties: Double-Tag, CEM Channel]{Per-sample rate uncertainty (\\% up\/down). Double-tag, CEM channel.}\\label{tab:2tag_cem}\n \\end{table}\n\n \\begin{table}\n \\begin{scriptsize}\n\\begin{tabular}{ c c c c c c c c c c c c } \\toprule\n & nonW & $t\\bar{t}$ & s-top $s$ & s-top $t$ & Zjets & W+LF & Wbb & Wcc & Wcj & WW & WZ\/ZZ \\\\\\midrule\n$\\texttt{ XS\\_ttbar}$ & $0\/0$ & $10\/-10$ & $10\/-10$ & $10\/-10$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFBTAGSF}$ & $0\/0$ & $11\/-11$ & $11\/-11$ & $12\/-12$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $24\/-24$ & $11\/-11$ \\\\\n$\\texttt{ CDFLUMI}$ & $0\/0$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $6\/-6$ & $6\/-6$ \\\\\n$\\texttt{ LEPACC\\_PHX}$ & $0\/0$ & $2\/-2$ & $2\/-2$ & $2\/-2$ & $2\/-2$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $2\/-2$ & $2\/-2$ \\\\\n$\\texttt{ QCD\\_PHX}$ & $40\/-40$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFVHF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $30\/-30$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFWC}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFMISTAG}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $21\/-21$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ZJETS}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $45\/-45$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ISRFSRPDF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $4\/-4$ & $4\/-4$ \\\\\n$\\texttt{ CDFJES}$ & $0\/0$ & $-7\/11$ & $-5\/-2$ & $-9\/0$ & $1\/-5$ & $3\/-2$ & $-6\/5$ & $-17\/8$ & $-35\/32$ & $-50\/-50$ & $-4\/-4$ \\\\\n$\\texttt{ CDFQ2}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $5\/0$ & $6\/-7$ & $-0\/35$ & $-41\/9$ & $0\/0$ & $0\/0$ \\\\\n\\hline\n$\\texttt{TOT}$ & $40\/40$ & $17\/20$ & $17\/16$ & $19\/17$ & $45\/46$ & $22\/21$ & $31\/31$ & $34\/47$ & $62\/45$ & $56\/56$ & $14\/14$ \\\\\n\\hline\n\\end{tabular}\n \\end{scriptsize}\n\\caption[Rate Uncertainties: Double-Tag, PHX Channel]{Per-sample rate uncertainty (\\% up\/down). Double-tag, PHX channel.}\\label{tab:2tag_phx}\n \\end{table}\n\n \\begin{table}\n \\begin{scriptsize}\n\\begin{tabular}{ c c c c c c c c c c c c } \\toprule\n & nonW & $t\\bar{t}$ & s-top $s$ & s-top $t$ & Zjets & W+LF & Wbb & Wcc & Wcj & WW & WZ\/ZZ \\\\\\midrule\n$\\texttt{ XS\\_ttbar}$ & $0\/0$ & $10\/-10$ & $10\/-10$ & $10\/-10$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFBTAGSF}$ & $0\/0$ & $11\/-11$ & $11\/-11$ & $12\/-12$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $24\/-24$ & $11\/-11$ \\\\\n$\\texttt{ CDFLUMI}$ & $0\/0$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $6\/-6$ & $6\/-6$ \\\\\n$\\texttt{ LEPACC\\_MU}$ & $0\/0$ & $1\/-1$ & $1\/-1$ & $1\/-1$ & $1\/-1$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $1\/-1$ & $1\/-1$ \\\\\n$\\texttt{ QCD\\_MU}$ & $40\/-40$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFVHF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $30\/-30$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFWC}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFMISTAG}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $21\/-21$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ZJETS}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $45\/-45$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ISRFSRPDF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $4\/-4$ & $4\/-4$ \\\\\n$\\texttt{ CDFJES}$ & $0\/0$ & $-7\/9$ & $-3\/2$ & $-2\/0$ & $7\/2$ & $-0\/0$ & $-10\/9$ & $-6\/10$ & $-14\/7$ & $-1\/-6$ & $-2\/-3$ \\\\\n$\\texttt{ CDFQ2}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $-2\/-5$ & $14\/10$ & $-6\/-8$ & $-8\/-11$ & $0\/0$ & $0\/0$ \\\\\n\\hline\n$\\texttt{TOT}$ & $40\/40$ & $18\/19$ & $16\/16$ & $17\/17$ & $46\/45$ & $21\/22$ & $35\/33$ & $31\/33$ & $34\/33$ & $25\/26$ & $13\/13$ \\\\\n\\hline\n\\end{tabular}\n \\end{scriptsize}\n\\caption[Rate Uncertainties: Double-Tag, CMUP $+$ CMX Channel]{Per-sample rate uncertainty (\\% up\/down). Double-tag, CMUP $+$ CMX channel.}\\label{tab:2tag_mu}\n \\end{table}\n\n\n \\begin{table}\n \\begin{scriptsize}\n\\begin{tabular}{ c c c c c c c c c c c c } \\toprule\n & nonW & $t\\bar{t}$ & s-top $s$ & s-top $t$ & Zjets & W+LF & Wbb & Wcc & Wcj & WW & WZ\/ZZ \\\\\\midrule\n$\\texttt{ XS\\_ttbar}$ & $0\/0$ & $10\/-10$ & $10\/-10$ & $10\/-10$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFBTAGSF}$ & $0\/0$ & $11\/-11$ & $11\/-11$ & $12\/-12$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $24\/-24$ & $11\/-11$ \\\\\n$\\texttt{ CDFLUMI}$ & $0\/0$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $6\/-6$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $6\/-6$ & $6\/-6$ \\\\\n$\\texttt{ LEPACC\\_EMC}$ & $0\/0$ & $5\/-5$ & $5\/-5$ & $5\/-5$ & $5\/-5$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $5\/-5$ & $5\/-5$ \\\\\n$\\texttt{ QCD\\_EMC}$ & $40\/-40$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFVHF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $30\/-30$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFWC}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $30\/-30$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ CDFMISTAG}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $21\/-21$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ZJETS}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $45\/-45$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ \\\\\n$\\texttt{ ISRFSRPDF}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $4\/-4$ & $4\/-4$ \\\\\n$\\texttt{ CDFJES}$ & $0\/0$ & $-7\/9$ & $-0\/-0$ & $1\/-5$ & $5\/-3$ & $-1\/1$ & $-10\/11$ & $-5\/19$ & $-7\/6$ & $-10\/-14$ & $8\/-8$ \\\\\n$\\texttt{ CDFQ2}$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $0\/0$ & $-3\/-1$ & $17\/5$ & $15\/0$ & $27\/20$ & $0\/0$ & $0\/0$ \\\\\n\\hline\n$\\texttt{TOT}$ & $40\/40$ & $18\/19$ & $17\/17$ & $18\/18$ & $46\/46$ & $21\/21$ & $36\/32$ & $34\/36$ & $41\/37$ & $27\/29$ & $16\/16$ \\\\\n\\hline\n\\end{tabular}\n \\end{scriptsize}\n\\caption[Rate Uncertainties: Double-Tag, EMC Channel]{Per-sample rate uncertainty (\\% up\/down). Double-tag, EMC channel.}\\label{tab:2tag_emc}\n \\end{table}\n\n\\renewcommand{\\tabcolsep}{6pt}\n\n\n\\chapter{Kinematic Distribution of the Signal Regions}\n\nFigures from~\\ref{fig:1svt_jet} to~\\ref{fig:2svt_dphi} show several kinematic variables for the single-tag and double-tag signal regions after the composition of all the background and all the lepton categories. The normalization of each background has been changed to match the {\\em best fit} values returned by statistical analysis of the likelihood described in by Equation~\\ref{eq:likelihood_full}, in Chapter~\\ref{chap:StatRes}. For each variable we show:\n\\begin{itemize}\n\\item the total background prediction overlaid to the selected data.\n\\item The reduced $\\chi^2$ of the prediction against the data distribution. As we use the best-fit normalization is used for the backgrounds, no rate uncertainty is applied in the $\\chi^2$ evaluation, however MC statistical uncertainty is accounted.\n\\item The probability of the Kolmogorov-Smirnov test derived from the predicted and observed shapes.\n\\item Background subtracted data histogram.\n\\item JES and $Q^2$ shape variations added in quadrature.\n\\end{itemize}\n\n\\begin{sidewaysfigure}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_J1Et}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_J2Et}\\\\\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_J1Eta}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_J2Eta}\n \\caption[Single Tag Region Jet Kinematic Variables]{Variables relative to the jet kinematic for all the lepton categories combined in the single-tag signal region. Jet 1 $E_T$ (top left), jet 2 $E_T$ (top right), jet 1 $\\eta$ (bottom left), jet 2 $\\eta$ (bottom left).}\\label{fig:1svt_jet}\n\\end{sidewaysfigure}\n\n\\begin{sidewaysfigure}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_LepPt}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_LepEta}\\\\\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_met}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_WMt}\n \\caption[Single Tag Region Lepton Kinematic Variables]{Variables relative to the lepton kinematic for all the lepton categories combined in the single-tag signal region. Lepton $P_T$ (top left), lepton $\\eta$ (top right), \\mbox{${\\not}{E_T}$} (bottom left), $M_T^W$ (bottom left).}\\label{fig:1svt_lep}\n\\end{sidewaysfigure}\n\n\\begin{sidewaysfigure}\n\\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_dPhiMetJet1}\n\\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_dRLepJ1}\n\\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_dRJ1J2}\n\\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVT_CEMCMUPCMXISOTRKPHX_wPt}\\\\\n\\caption[Single Tag Region Angular Kinematic Variables]{Angular variables for all the lepton categories combined in the single-tag signal region. $\\Delta\\phi($\\mbox{${\\not}{E_T}$} $,jet1)$ (top left), $\\Delta R(Lep,jet1)$ (top right), $\\Delta R(jet1,jet2)$, (bottom left), $P_T^W$ (bottom left).}\\label{fig:1svt_dphi}\n\\end{sidewaysfigure}\n\n\n\\begin{sidewaysfigure}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_J1Et}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_J2Et}\\\\\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_J1Eta}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_J2Eta}\n \\caption[Double Tag Region Jet Kinematic Variables]{Variables relative to the jet kinematic for all the lepton categories combined in the double tag-signal region. Jet 1 $E_T$ (top left), jet 2 $E_T$ (top right), jet 1 $\\eta$ (bottom left), jet 2 $\\eta$ (bottom left).}\\label{fig:2svt_jet}\n\\end{sidewaysfigure}\n\n\\begin{sidewaysfigure}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_LepPt}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_LepEta}\\\\\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_met}\n \\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_WMt}\n \\caption[Double Tag Region Lepton Kinematic Variables]{Variables relative to the lepton kinematic for all the lepton categories combined in the double tag-signal region. Lepton $P_T$ (top left), lepton $\\eta$ (top right), \\mbox{${\\not}{E_T}$} (bottom left), $M_T^W$ (bottom left).}\\label{fig:2svt_lep}\n\\end{sidewaysfigure}\n\n\\begin{sidewaysfigure}\n\\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_dPhiMetJet1}\n\\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_dRLepJ1}\n\\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_dRJ1J2}\n\\includegraphics[width=0.49\\textwidth]{AppTagKin\/Figs\/WZ_2JET_SVTSVT_CEMCMUPCMXISOTRKPHX_wPt}\\\\\n\\caption[Double Tag Region Angular Kinematic Variables]{Angular variables for all the lepton categories combined in the double-tag signal region. $\\Delta\\phi($\\mbox{${\\not}{E_T}$} $,jet1)$ (top left), $\\Delta R(Lep,jet1)$ (top right), $\\Delta R(jet1,jet2)$, (bottom left), $P_T^W$ (bottom left).}\\label{fig:2svt_dphi}\n\\end{sidewaysfigure}\n\n\n\\chapter{WHAM: WH Analysis Modules}\\label{chap:AppWHAM}\n\nHigh energy physics data analysis is done with the help of complex software frameworks which allow to manage the huge amount of information collected by the detector. \nThe analysis framework which was used for this analysis is named $WH$ {\\em Analysis Modules} or WHAM and I was one of the main developers of the package. \n\nThe aim of the software is a reliable event selection and background estimate in the $\\ell\\nu+HF$ channel, with the possibility to easily implement new features and studies that can improve the Higgs search.\nWHAM plays a relevant role in the CDF low-mass Higgs boson search, especially in the $WH\\to \\ell\\nu+b\\bar{b}$ channel but also in other contexts, like $t\\bar{t}H$ and lately also $ZH$ searches. \n\nThe analysis package tries to incorporate the CDF knowledge about the $\\ell\\nu+HF$ channel, most of it coming from the top~\\cite{top_lj2005} and single-top~\\cite{single_top} analyses. An effort was also made in the direction of code modularization and analysis customization with option loading at run-time.\n\nA detailed explanation of the package is beyond the goal of this thesis, however a general overview of the package structure and functionality is given. More information is available in the CDF internal pages~\\cite{wham_twiki} although a comprehensive documentation is not yet available.\n\n\n\\section{Package Structure}\n\nThe package is organized in a folder structure organized according to the purpose of each of the sub-elements. The {\\em first level} directories are:\n\\begin{description}\n\\item[Setup:] contains the scripts needed to setup the analysis environment, both first installation and every day use, the references to all the external tools and any patche that needs to be applied.\n\\item[Documentation:] contains all the internal and public documentation. It is easily accessible and customizable by all the analysis group collaborators.\n\\item[Inputs:] contains the database files of the MC samples, the parametrization of the triggers and $SF$s, the re-weighting templates, the option configuration files. Basically every input to the analysis elaboration is here, except the data and the MC ntuples themselves.\n\\item[Commands:] this is more an utilities repository, it contains scripts to run the analysis on the CDF Grid for parallel computing (named CAF~\\cite{caf_Sfiligoi2007}), plus a wide set of macros and scripts used for single-sample studies, text file processing or small data-handling tasks. \n\\item[Results:] contains all the information elaborated by the rest of the analysis packages. This includes pre-processed data and MC samples as well as the final templates obtained after the complete elaboration. Several commands expect to find the input files here.\n\\item[Modules:] this is the core of the analysis package. It contains the C++ code used for the selection, the background estimate and the final production of the templates used for the statistical analysis and the validation of the kinematic distributions. Next section will describe it in more details.\n\\end{description}\n\n\\subsection{WHAM Modules}\n\nThe core of the WHAM analysis framework is the Modules directory. Here each functional step of the analysis is classified in a {\\em module}, i.e. a self consistent C++ class built with standardized structure to allow straightforward compilation and testing. \n\nThe modules are of two kinds: {\\em functional modules} and {\\em construction modules}. The formers are in charge of actually perform an operation, for example the event selection or the drawing of stacked histograms. The second kind of modules are the sub-components used by the first, for example the selection code needs to known the format of the input and output data as well as the definition of the lepton-object or jet-object.\n\nThe level of abstraction offered by the building modules proved to be extremely powerful. Two minimal examples (on which I contributed) that revealed to be extremely useful are: the handling of the configuration options and of the $b$-tag efficiency estimate. \nFor the first, I implemented, using the \\verb'libconfig' library~\\cite{libconfig}, a text file reading utility that allows a single location definition and the run-time loading of all the options needed by the functional modules. For the second, it was necessary to identify the minimal amount of information needed to define a $b$-tagging algorithm. The only two values needed\\footnote{Functional dependencies and correlations should be already taken into account.} are: $SF_{Tag}$ and $p_{Mistag}^j$ (see Section~\\ref{sec:tag-sel}). With this information, it was possible to develop a single algorithm for the combination of any number of different $b$-tagging algorithms, for any required tag and jet multiplicity. The code works iteratively on the jets of the event requiring the definition of a $b$-tag in priority order, defined by the user.\n\n\nThe functional modules are four and, in the directory structure, are identified by the \\verb'process' prefix:\n\\begin{description}\n\\item[Sample Selection:] access to the production ntuples (see Section~\\ref{sec:data-struc}), as well as lepton and jet selections, are performed here. The result is a small size ntuple, the \\texttt{EvTree}, containing the 4-vectors of the identified particles plus all the relevant information needed in the next steps. The \\texttt{EvTree} is defined by a class with complex {\\em methods} working on the simple stored variables: this allows a huge saving of disk space and computing time at selection stage. Furthermore the portability of the ntuple class ease the reproducibility of the same algorithms and allows faster checks.\n\\item[Sample Pre-processing:] here the complete selection is applied to the \\texttt{EvTree}'s and the pre-processing with more analysis-specific algorithms is performed. For example, MC samples are scaled to the expected yield with the application of the latest available $SF$'s and trigger efficiencies. Also, the multivariate discriminants are evaluated here. The pre-processing has the possibility to be interfaced to other ntuples than the \\texttt{EvTree}.\n\\item[Background Estimate:] this is the last step in the analysis of the $\\ell\\nu+HF$ channel. The background estimate described in Chapter~\\ref{chap:bkg} is applied here and the templates of the different signals, backgrounds and data samples are stored in a ROOT file with the derived normalizations. In general any other background estimate method can be plugged at the end of the analysis chain but for the moment only the one described in Chapter~\\ref{chap:bkg} is available.\n\\item[Result Display:] this step completes the analysis in the sense that it allows the comparison of the final estimate with the observed data distributions. Histograms are produced together with statistical indicators of the shape agreement: $\\chi^2$, Kolmogorov-Smirnov tests, systematic overlay, background subtracted plots. Tables and histograms can appear in several formats, from html pages to simple eps files.\n\\end{description}\nThe very last step in each analysis, the statistical interpretation of the results, is implemented in a different software tool~\\cite{mclimit} that is interfaced with the templates produced by the WHAM background estimate.\n\n\n\\section{Relevant Results}\n\nBeyond the analysis presented here, a wide range of other analyses exploit the WHAM package. Between the most relevant: the new single-top cross section measuremens~\\cite{single_top_benwu_10793}, the $t\\bar{t}H$~\\cite{tth_10801} search in the lepton plus jets channel and, lately, also top-properties~\\cite{topR_10887} and SUSY searches are exploiting the package.\n\nFigure~\\ref{fig:WH_improve}, probably, shows the most striking result: the improvement of the $WH$ search sensitivity (7.5~fb$^{-1}$ and 9.4~fb$^{-1}$ versions~\\cite{wh75_note10596, wh94_note10796}). The several improvements produced for the Higgs search were readily implemented and tested thanks to the backbone of a reilable framework. \n\nNevertheless the shrinking of the CDF collaboration, the final sensitivity to the Higgs boson exceeded the best expectations.\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.85\\textwidth,height=0.75\\textwidth]{AppWHAM\/Figs\/WH_Projection_new}\n\\caption[Improvements to the CDF $WH$ Search Sensitivity]{Improvements to the CDF $WH\\to \\ell\\nu+ b\\bar{b}$ search sensitivity for a Higgs boson of $m_H=115$~GeV$\/c^{2}$. Both the 2011~\\cite{wh75_note10596} and 2012~\\cite{wh94_note10796} results are mainly obtained within the WHAM analysis framework. The sensitivity reached improves more than 125\\% over the first analysis of the same channel.}\\label{fig:WH_improve}\n\\end{center}\n\\end{figure}\n\n\n\n\\chapter{Background Estimate}\\label{chap:bkg}\n\nIn this analysis an accurate background evaluation is crucial since a low signal over background ratio is expected and several components contribute to the final $M_{Inv}$ line shape used for the signal extraction. \n\nBackground estimate should be as independent as possible from the signal regions, the single and double-tagged $\\ell\\nu+2 $~jets samples. One possible approach is the simulation of background events with Monte Carlo generators, another approach is the extrapolation of background components from data control regions. Both methods present advantages and disadvantages, for example a MC simulation is completely under control but can be limited by the knowledge of the underlying physics phenomena parametrization. Data-driven techniques, on the other hand, are based on the assumption that a background in the control region can be simply extrapolated to the signal region. This brings some approximations, eventually covered by systematic uncertainties. \n\nIn our case we use a mixture of MC and data-driven estimates depending on the different background component being addressed.\n\nThe complete machinery used for the $\\ell\\nu+HF$ jets background estimate is named {\\em Method II}\\footnote{The name derives from a methodology developed for the lepton plus jets $t\\bar{t}$ decay channel. Method II appeared after the {\\em Method I}, the data-driven background estimate used in the top discovery~\\cite{cdf_top_evidence}.} and it was extensively studied at CDF, first on $t\\bar{t}$~\\cite{top_lj2005} analyses, and later in single-top~\\cite{single_top} and $WH$~\\cite{wh75_note10596, wh94_note10796} analyses. Together with the MC, two data control samples play a major role in this kind of background estimate: the {\\em pretag} control sample, without any $b$-tag requirement, and the fake-$W$ enriched sample, without any multi-jet rejection requirement.\n\nThe background evaluation can be divided in four categories: \n\\begin{description}\n\\item[Electroweak and Top:] sometimes dubbed $EWK$, these are backgrounds due to contributions of well known physical processes: $t\\bar{t}$, single-top and $Z+ $~jets production. For these processes, described in detail in Section~\\ref{sec:mc_bkg}, we completely rely on MC simulation for the line-shape and theoretical prediction or previous measurements for the normalization. \n\\item[$W+$ Heavy Flavors:] also named $W+HF$, is the part of the $W+$~jets sample where the $W$ is produced in association with $HF$ quarks. It is classified as {\\em irreducible} because it has the same final state ($W+c\/b$ quarks) of the $diboson\\to \\ell\\nu+HF$ signal. The line shape is derived by a LO $W+n$~partons MC while the normalization is obtained from the {\\em pretag} data control region and from MC derived $W+HF$ fractions. Section~\\ref{sec:whf} describes in detail the procedure.\n\\item[$W+$ Light Flavors:] also named {\\em mistags} or $W+LF$. These are events where the $W$ is produced in association with one or more Light Flavor ($LF$) jets. This component is drastically reduced by the requirement of one or more $\\mathtt{SecVtx}$ tagged jets. Remaining events are due to mis-measured jet tracks and long living $LF$ hadrons. The background evaluation, explained in Section~\\ref{sec:mistag}, is obtained using LO $W+ n$~partons MC where appropriate reweighting and normalization are derived from a data control sample of multi-jet events.\n\\item[Fake-$W$:] sometimes also named multi-jet or non-$W$ background, is composed by events without a real $W$ passing the selection. Those are due to multi-jet QCD production where the jets fake the charged lepton and the \\mbox{${\\not}{E_T}$} of the neutrino. Section~\\ref{sec:qcd} describes the normalization and line-shape evaluation procedure, both completely data-driven.\n\\end{description}\nThe normalization of the four categories proceeds through sequential steps: the $EWK$ contribution is evaluated first, then $W+$~jets normalization is extracted from the pretag control region so that the $W+HF$ fraction can be scaled to the single and double tag signal regions.\n\nNext step is the estimate, from data, of $W+LF$ contribution and, finally, the {\\em tagged} non-$W$ contribution is evaluated before the application of the SVM requirement.\n\n\n\\section{Electroweak and Top Processes}\\label{sec:mc_bkg}\n\nThis category contains all the processes whose shape is evaluated directly from MC samples and whose normalization is obtained from theory prediction or experimental measurements. Figure~\\ref{fig:fey_ewk} shows some of the tree-level Feynman diagrams contributing to this background set. \n\nThe components of $EWK$ backgrounds are: \n\\begin{description}\n\\item[$Z\/\\gamma$ + jets:] Drell-Yan and $Z\\to \\ell\\ell$ events can be produced in association with jets emitted by QCD radiation. If one lepton is misidentified, they enter as a background in our selection. The simulation of $Z\/\\gamma+n$ partons is available only at LO level (with \\texttt{ALPGEN}~\\cite{alpgen} Matrix Element generator) therefore the process is evaluated with the same procedure of $W+$~jets background (see Section~\\ref{sec:whf}). Luckily, the small acceptance and the possibility to use a cross section measured at CDF~\\cite{z_jets_cx} allows us to evaluate the $Z\/\\gamma+$jets background directly from MC. The $Z+$jets cross section measurement refers only to the on-shell $Z$ contribution, therefore a further scaling factor is derived by comparing the \\texttt{ALPGEN} LO predictions for the on-shell and the inclusive $Z\/\\gamma+$ jets MC sample. \n\n\\item[Top pair:] $t\\bar{t}$ production and decay in the lepton plus jets ($\\ell\\nu+b\\bar{b}+2$ jets) and di-lepton ($2\\ell+2\\nu+b\\bar{b}$) channels can enter in our selection if one or more of the final state objects is misidentified. Although the probability is small, the background becomes sizable in the $HF$-tagged sample because of the presence of two $b$-quarks. \n\\item[Single-top:] $t$ and $s$ channel production have cross sections comparable to the diboson one and the semi-leptonic decay channel produces the same final state we are interested in. On the other hand, the signal to background discrimination is performed by the $M_{Inv}(jet1,jet2)$ where the single-top has no peak.\n\\end{description}\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.25\\textwidth]{Bkg\/Figs\/zjet}\n\\includegraphics[width=0.7\\textwidth]{Bkg\/Figs\/ttbar}\\\\\n\\includegraphics[width=0.7\\textwidth]{Bkg\/Figs\/single-top_s_t}\n\\caption[Some Tree-Level Diagrams Contributing to the $EWK$ Backgrounds]{Examples of tree-level Feynman diagrams contributing to the $EWK$ background set. They are $Z+$jets production (top left in association with a $b\\bar{b}$ quark pair), $t\\bar{t}$ production and decay in semi-leptonic (top center) and di-leptonic channels (top right), single-top production in the $s$ (bottom left) and $t$ (bottom right) channels and semi-leptonic decay.}\\label{fig:fey_ewk}\n\\end{center}\n\\end{figure}\n\nThe normalization, $N_{l,k}^{EKW}$, of each $EWK$ process is obtained for each lepton category, $l$, in the pretag and $1,2$-Tag regions, $k$, by using Equation~\\ref{eq:mc_estimate}. A summary of the production cross sections and the MC generator information are reported in Table~\\ref{tab:mc_ewk}.\n\n \\begin{table}[h]\n \\begin{center}\n \\begin{tabular}{cccc}\n \\toprule\n Sample & $\\sigma$ (pb)& Source & Initial Events \\\\\n \\midrule\n $Z\/\\gamma+$jets & 787.4 $\\pm$ 85 & Measured~\\cite{z_jets_cx} & Several MCs \\\\ \n $t\\bar{t}$ & 7.04 $\\pm$ 0.49 & NNLO~\\cite{ttbar_cx} & 6.7 M\\\\ \n single-top, $s$ & 1.017 $\\pm$ 0.065 & NNLO~\\cite{stops_kidonakis} & 3.5 M\\\\ \n single-top, $t$ & 2.04 $\\pm$ 0.18 & NNLO~\\cite{stopt_kidonakis} & 3.5 M\\\\ \n \\bottomrule\n \\end{tabular}\n \\caption[MC Information for the $EWK$ Backgrounds]{MC information for backgrounds of the $EWK$ category, i.e. the ones estimated from simulation. $Z+$ jets cross section~\\cite{z_jets_cx} refers to the on-shell $Z$ contribution and \\texttt{ALPGEN} LO MC prediction is used to derive a correction factor. \\texttt{HERWIG} plus \\texttt{PYTHIA} PS are used for $t\\bar{t}$, \\texttt{POWHEG} plus \\texttt{PYTHIA} PS are used for single-top $s$ and $t$ channels. The generated top mass is $m_t=172.5$~GeV$\/c^{2}$.}\\label{tab:mc_ewk}\n \\end{center}\n \\end{table}\n\nAs the three MC derived background processes and the diboson signal are estimated in the same way, the label $EWK$ will be used to indicate generically all of them:\n\\begin{equation}\n \\sum\\limits_{EWK} N_{l,k}^{EWK} = N_{l,k}^{Zjets} + N_{l,k}^{t\\bar{t}} + N_{l,k}^{s-top,s} + N_{l,k}^{s-top,t} + \\sum\\limits_{diboson} N_{l,k}^{diboson}\\mathrm{.}\n\\end{equation}\n\n\n\n\\section{$W$ plus Heavy Flavor Background}\\label{sec:whf}\n\nThe associated $W+HF$ production (Figure~\\ref{fig:fey_wjets} shows few examples of tree-level diagram contributions) is the most significant contribution to the total background in the single and double tag signal samples. For its estimate we rely partially on MC and partially on data. \n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.9\\textwidth]{Bkg\/Figs\/wjet}\n\\caption[Some Tree Level Diagrams Contributing to $W+HF$ and $W+LF$ Backgrounds]{Examples of tree-level Feynman diagrams contributing to the $W+HF$ and the $W+LF$ backgrounds. NLO diagrams also play a relevant role in the $W+$~QCD radiation emission.}\\label{fig:fey_wjets}\n\\end{center}\n\\end{figure}\n\nNormalization and kinematic distributions of the $W+$~jets sample, and of the $W+HF$ sub-sample, present several theoretical difficulties. NLO MC simulations became available lately for $W+b\\bar{b}$ processes~\\cite{wbb_powheg} but large higher order contributions enhance the theoretical uncertainty. Furthermore the prediction of the complete $HF$ spectrum ($W+c\\bar{c}$, $W+b\\bar{b}$, $W+c$) is beyond the current capabilities. \n\nThe two aspects of the problem can be solved by factorizing the evaluation of the kinematic properties from the estimate of the normalization:\n\\begin{itemize}\n\\item a solution to estimate the kinematic properties and the composition of $W+$~jets QCD emission was proposed in~\\cite{first_lo_match} in 2001: the appropriate matching scheme of a LO calculation with a properly tuned Parton Shower (PS) simulation can reproduce the experimental results. Section~\\ref{sec:alpgen_comp} describes in detail our use of the \\texttt{ALPGEN} LO generator, interfaced with \\texttt{PHYTIA} PS. $W+HF$ relative components or $HF$ fractions, $f^{HF}$, can be extracted with this method.\n\\item the only free parameter remaining is, now, the $W+$~jets normalization. We derive it completely from data using a maximum likelihood fit on the SVM output distribution: this variable, described in Appendix~\\ref{chap:AppSvm}, has a high discriminating power between non-$W$ and real $W$ events. Therefore, as all the other $EWK$ components are known, both $N^{nonW}_{pretag}$ and $N^{Wjets}_{pretag}$ can be extracted from the fit. Section~\\ref{sec:wjets_pretag} describes the procedure.\n\\end{itemize}\n\n\n\n\\subsection{$W +$ Jets Spectrum Composition}\\label{sec:alpgen_comp}\n\n$W+$~jets composition and kinematic properties are evaluated from MC. We employ a large set of $W+n$~partons \\texttt{ALPGEN}~\\cite{alpgen} LO MC that includes $HF$ production in the ME calculation and is interfaced with \\verb'PYTHIA'~\\cite{pythia} PS:\n\\begin{itemize}\n\\item[-]$W$ plus $0$, $1$, $2$, $3$, $\\ge 4$ partons (for $W+ n$~generic jets description);\n\\item[-]$W+b\\bar{b}$ plus $0$, $1$, $2$ partons;\n\\item[-]$W+c\\bar{c}$ plus $0$, $1$, $2$ partons;\n\\item[-]$W+c$ (or $W+\\bar{c}$) plus $0$, $1$, $2$, $3$ partons.\n\\end{itemize}\nIn total we used ninety samples: fifteen to be multiplied by the three leptonic $W$ decay modes ($W\\to e,\\mu,\\tau +\\nu_e,\\nu_{\\mu}\\nu_{\\tau}$) and by another factor two for {\\em low luminosity} and {\\em high luminosity} running periods. Table~\\ref{tab:samples} summarizes the total of ninety samples used along with the respective LO cross sections and the approximate number of events initially generated.\n\\begin{table}\\begin{center}\n\\begin{tabular}{ccc}\n\\toprule\nSample& $\\sigma$ (pb)& Initial Events\\\\\n\\midrule\n$W+0p$& 1810& 7 M\\\\\n$W+1p$& 225 & 7 M\\\\\n$W+2p$& 35.3& 1.4 M\\\\\n$W+3p$& 5.59& 1.4 M\\\\\n$W+4p$& 1.03& 0.5 M\\\\\n$W+b\\bar{b}+0p$& 2.98& 2.1 M\\\\\n$W+b\\bar{b}+1p$& 0.888& 2.1 M\\\\\n$W+b\\bar{b}+2p$& 0.287& 2.1 M\\\\\n$W+c\\bar{c}+0p$& 5& 2.8 M\\\\\n$W+c\\bar{c}+1p$& 1.79& 2.8 M\\\\\n$W+c\\bar{c}+2p$& 0.628& 2.8 M\\\\\n$W+c+0p$& 17.1& 2.8 M\\\\\n$W+c+1p$& 3.39& 2.8 M\\\\\n$W+c+2p$& 0.507& 2.8 M\\\\\n$W+c+3p$& 0.083& 2.8 M\\\\\n\\bottomrule\n\\end{tabular}\n\\caption[ALPGEN $W + $ Partons MC Samples]{\\texttt{ALPGEN} LO plus \\texttt{PYTHIA} (PS) MC samples used to estimate the various contributions to the $W+HF$ and $W+LF$ backgrounds. The second column lists the LO cross section of the samples, the third column shows the approximate number of generated events. Different MC samples are used for the three lepton flavors ($e, \\mu, \\tau$) of the $W$ decay.}\\label{tab:samples}\n\\end{center}\\end{table}\n\nThe basic composition algorithm is simple: the acceptance of each sample, for a given jet multiplicity and for each lepton category, is weighted by its own LO production cross section, then a normalization factor equal to the sum of all the samples is applied. The key idea is that, if the PS matching is well tuned, the higher order corrections to the LO $W+n$~partons simulation, should simplify when the normalization factor is applied. \n\nThe procedure is more complex in our case because we want an accurate estimate of the $W+HF$ spectrum: we need to derive the $HF$ fractions, $f^{HF}$, from the $W+n$~partons samples and relate them to the physical observables of $W+b\\bar{b}\/c\\bar{c}\/c$ jets with one or two \\texttt{SecVtx} $HF$-tags.\n\nThe $W+HF$ production is evaluated explicitly at ME level in some of the $W+n$ partons samples, however a certain amount of $HF$ quarks is also produced by the PS algorithm in all the samples. Some care must be used when combining all the samples because the double counting of a process will degrades the relative weights.\n\nThe {\\em overlap removal} technique used in this work was previously developed by the CDF collaboration and it is documented in~\\cite{hf_8765}. It consists in the selection of $HF$ production from the PS {\\em or} from the ME on the base of detector level reconstructed jets. We explicitly veto events where the ME heavy flavor quarks wind up in the same jet, and we also remove events where the $HF$ quark pair from the shower is divided in two jets. In both cases, the distinction is made with simple kinematic cuts that define a parton to be {\\em in a jet} if:\n\\begin{equation}\n \\Delta R(parton, jet)<0.4\\quad\\textrm{and}\\quad E_T^{jet,cor}> 12~\\textrm{GeV}\n\\end{equation}\nThe reason for this choice is that PS and ME predict different $\\Delta R$ (separation) between heavy quarks: PS model is tuned on the more collinear gluon splitting quark pairs but fails in the limit of large opening angles, while ME generation works better in latter regime. Figure~\\ref{fig:overlap_rem_spectrum} shows the generator level $\\Delta R$ distribution of $b\\bar{b}$ and $c\\bar{c}$ pair production after the composition of the $W+n$~partons samples and the overlap removal procedure. The $\\Delta R $ from PS-only is also reported in Figure~\\ref{fig:overlap_rem_spectrum} showing disagreement at large $\\Delta R $ with respect to the $HF$ generated with ME. The smooth transition between the different $W+n$~partons contributions indicates that the overlap removal scheme is working properly.\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/cc_dr}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/bb_dr}\n\\caption[Generator Level $\\Delta R$ Distributions of $b\\bar{b}$ and $c\\bar{c}$ ]{Generator level $\\Delta R$ distributions of $b\\bar{b}$ (left) and $c\\bar{c}$ (right) pair production after the composition of the $W+n$~partons samples and the overlap removal procedure~\\cite{hf_8765}. The $\\Delta R $ from PS-only is also reported showing disagreement at large $\\Delta R $ w.r.t. the $HF$ generated with ME.}\\label{fig:overlap_rem_spectrum}\n\\end{center}\n\\end{figure}\n\nThe complete procedure for the determination of the $f^{HF}$'s is the following:\n\\begin{itemize}\n\\item each MC sample is assigned a weight:\n \\begin{equation}\\label{eq:weight}\n w^{m}=\\frac{\\sigma^m}{N_{z_0}^m} \\mathrm{,}\n \\end{equation}\n where $m$ varies on all the different samples of Table~\\ref{tab:samples}, $\\sigma^m$ is the cross section and $N_{z_0}^m$ is the total number of generated events with $|z_0^{vtx}|<60$~cm.\n\\item After pretag selection, the denominators of the $HF$ fractions are calculated as:\n \\begin{equation}\n Den=\\sum_{m=1}^{samples} w^{m}N^m\\mathrm{,}\n \\end{equation}\n where $N^m$ are the {\\em unique} events (i.e. not vetoed by the overlap removal) falling in each jet-bin for sample $m$.\n\\item Each selected jet is classified as $HF$ jet or not: $HF$ jets have a $b$ or $c$ parton (at generation level) {\\em match}, i.e. laying inside the cone (\\mbox{$\\Delta R(jet,parton)<0.4$}). We distinguish four $HF$ categories depending on the number of matched jets and on the kind of $HF$ parton ($b$ or $c$). The categories are: $1B$, $2B$, $1C$ and $2C$ depending if we have 1 or $\\ge 2$, $b$ or $c$ quarks matched; no $b$ quarks should be matched in the case of $1,2C$ categories.\n\\item The numerators are defined by the sum of the events in each $HF$ category (with weights given by Equation~\\ref{eq:weight}) over all the samples, in each jet-bin:\n \\begin{equation}\n Num^{HF}=\\sum_{m=1}^{samples}w^m N^{m,HF}\\mathrm{.}\n \\end{equation}\n\\end{itemize}\nThe computation of the $f^{HF}$'s is now straightforward. Table~\\ref{tab:hf_frac} summarizes the final $HF$ fractions averaged on all the lepton categories.\n\\begin{table}\\begin{center}\n\\begin{tabular}{cccc}\n\\toprule\n\\multicolumn{4}{c}{Heavy Flavor Fractions}\\\\\n\\midrule\n$f^{1B}$ & $f^{2B}$ & $f^{1C}$ & $f^{2C}$ \\\\\n$0.024$ & $ 0.014$ & $ 0.114$ & $ 0.022$ \\\\ \n\\bottomrule\n\\end{tabular}\n\\caption[$HF$ fractions, $f^{HF}$, for Two Jets Selection]{$HF$ fractions, $f^{HF}$, for the two jets selection derived from \\texttt{ALPGEN} MC composition. They are classified in the different categories. $1B$: events with $1$ $b$-jet. $2B$: events with $\\ge 2$ $b$-jet. $1C$: events with $1$ $c$-jet and no $b$-jet. $2C$: events with $\\ge 2$ $c$-jet and no $b$-jet.}\\label{tab:hf_frac}\n\\end{center}\\end{table}\n\nHowever $f^{HF}$'s are not immediately usable to predict $W+HF$ background. They need an additional correction, $K$, to account for different $HF$ composition in data and in the LO ME simulation:\n\\begin{equation}\\label{eq:k}\nK^{HF}=\\frac{f^{HF}_{data}}{f^{HF}_{MC}}\\mathrm{,}\n\\end{equation}\nThe $HF$ calibration is carried out on a data sample not used in this measurement: we base our correction on the studies performed in~\\cite{single_top}. There, the $W+1$~jet data sample is used to derive a correction of \\mbox{$K^{c\\bar{c}\/b\\bar{b}}=1.4$} to the $W+c\\bar{c}\/b\\bar{b}$ processes, while no correction ($K^{c}=1.0$) is quoted for the $W+c$ process\\footnote{$W+c$ production is an electroweak production process measured at CDF in~\\cite{cdf_wc}.}. A large uncertainty of $30$\\% is associated to all these $K$ factors as the correction is extrapolated from the $W+1$~jet data sample to the $W\\ge 2$ jets data sample. A previous study~\\cite{k_8768}, on a multi-jet control sample, obtained a $K$ closer to $1$ and it is consistent within the uncertainty. Table~\\ref{tab:k-fac} shows a summary of the applied $K$ corrections.\n\\begin{table}\\begin{center}\n\\begin{tabular}{ccc}\n\\toprule\n\\multicolumn{3}{c}{Heavy Flavor $K-$Factor Corrections}\\\\\n\\midrule\n$K^{W+b\\bar{b}}$ & $K^{W+c\\bar{c}}$ & $K^{W+c}$ \\\\\n$1.4\\pm 0.42$ & $1.4\\pm 0.42$ & $1.0\\pm 0.33$ \\\\ \n\\bottomrule\n\\end{tabular}\n\\caption[$HF$ $K-$Factor Correction]{$K-$Factor corrections needed for the calibration of the $HF$ production in data. Studies performed in~\\cite{single_top}.}\\label{tab:k-fac}\n\\end{center}\\end{table}\n\nThe final estimate of $f^{HF}$ is:\n\\begin{equation}\nf^{HF}=K\\frac{Num^{HF}}{Den}\\mathrm{.}\n\\end{equation}\n\nThe $HF$ fractions can be used in the construction of the pretag sample shapes but are of marginal importance because of the overwhelming $LF$ contribution.\nOn the other hand $HF$ are fundamental in the tagged sample composition. To evaluate it we need the tagging efficiency of each $HF$ category ($\\epsilon^{HF}$): they are derived from the MC samples requiring a \\verb'SecVtx' tag, then we apply the $\\omega_{k-Tag}$ correction as described in Section~\\ref{sec:jetSel}. Tagging efficiencies for each $HF$ category are summarized in Table~\\ref{tab:hf_tag_ef}.\n\\begin{table}\\begin{center}\n\\begin{tabular}{ccccc}\n\\toprule\n\\multicolumn{5}{c}{Heavy Flavor Tagging Efficiency}\\\\\n\\midrule\n& $\\epsilon^{1B}$ & $\\epsilon^{2B}$ & $\\epsilon^{1C}$ & $\\epsilon^{2C}$ \\\\\n1-Tag & $ 22$\\% & $ 30$\\% & $6.7$\\% & $ 9.1$\\% \\\\ \n2-Tag & $ 0.4$\\% & $8.6$ & $ 0.06$\\% & $ 0.4$\\% \\\\ \n\\bottomrule\n\\end{tabular}\n\\caption[Heavy Flavor Tagging Efficiency $\\epsilon^{HF}$ ]{Heavy Flavor tagging efficiency $\\epsilon^{HF}$ in each category. $1B$: events with $1$ $b$-jet. $2B$: events with $\\ge 2$ $b$-jet. $1C$: events with $1$ $c$-jet. $2C$: events with $\\ge 2$ $c$-jet.}\\label{tab:hf_tag_ef}\n \\end{center}\n\\end{table}\n\n\nFinally we can combine all the pieces to obtain the {\\em final} $HF$ tagged contribution. For each lepton $l$ and tag category $k$, it is:\n\\begin{equation}\\label{eq:tag_hf}\nN^{HF}_{l,k}=N^{Wjets}_{l,k}\\cdot f^{HF} K^{HF}\\epsilon^{HF}_{k}\\mathrm{,}\n\\end{equation}\nWhere the $HF$ label distinguishes the number of jet-matched heavy flavor hadrons: $1,2B$, $1,2C$. The $1,2B$ fractions represent $W+b\\bar{b}$, $2C$ represents $W+c\\bar{c}$ prediction and $1C$ is related to $W+c$ production. The extraction of the $W+$~jets pretag normalization ($N^{Wjets}$) is discussed in the next Section.\n\n\n\n\\subsection{Pretag $W +$ Jets Normalization Estimate}\\label{sec:wjets_pretag}\n\nAccording to Equation~\\ref{eq:tag_hf}, the $W+$~jet normalization, $N^{Wjets}$, is the only missing piece in the complete $W+HF$ estimate. $N^{Wjets}$ is obtained solving:\n\\begin{equation}\\label{eq:pretag_comp}\nN^{Data}_{Pretag} = F^{nonW}_{Pretag} \\cdot N^{Data}_{Pretag} + F^{Wjets}_{Pretag} \\cdot N^{Data}_{Pretag} + \\sum_{EWK}N^{EWK}_{Pretag} \\mathrm{,}\n\\end{equation}\nwhere:\n\\begin{itemize}\n\\item $N^{EWK}$ is derived from Section~\\ref{sec:mc_bkg};\n\\item the fractions $F^{nonW}_{Pretag}$ and $F^{Wjets}_{Pretag}$ are unknown parameters obtained by a maximum likelihood fit in the pretag region. This provides a $W+$~jets estimate unbiased by the $b$-tag requirement;\n\\item the SVM output distribution is used in the fit before the application of the multi-jet rejection cut (Section~\\ref{sec:svm_sel}) so that the fake-$W$ fraction of the sample can be extrapolated from the sideband. \n\\end{itemize}\nWe build a binned likelihood function on the SVM distribution on the base of Equation~\\ref{eq:pretag_comp}. Events are Poisson distributed and a Gaussian constraint is imposed to the $N^{EWK}$ normalizations while $F^{nonW}_{Pretag}$ and $F^{Wjets}_{Pretag}$ are left free to vary.\n\nTemplates used in the fit are derived in different ways according to the background source: $EWK$ and $diboson$ SVM distributions are derived from MC, $W+$~jets template is derived from the $W+n$~partons MC composition described in previous Section. Non-$W$ template is obtained in data by reversing appropriate lepton identification cuts, a detailed description is given in Section~\\ref{sec:qcd}.\nTemplates are built for each lepton category but we decided to combine all the EMC algorithms in one template to reduce statistical fluctuations exploiting the homogeneous behaviour of the different components of EMC.\n\nThe validation of the SVM output distribution is another important step of the fit procedure. Most of the previous analyses~\\cite{single_top, wh94_note10796} used a fit to the \\mbox{${\\not}{E_T}$} distribution to derive the $W+$ jets fraction. Here a slightly different strategy is used as the SVM was trained to discriminate the $W+$~jets sample against multi-jet events. This physics interpretation was checked during the training process (see Appendix~\\ref{chap:AppSvm}). A simplified version of the fit (slightly different selection and composed only by non-$W$ and $W+$~jets samples) was performed at each training step to reject unreliable configurations. We based the rejection on two figures of merit: \n\\begin{itemize}\n\\item the $\\chi^2$ evaluation of post fit template\/data shapes; \n\\item the comparison of the non-$W$ and $W+$~jets expectation and fit estimate.\n\\end{itemize}\n\nThe final step is the actual fit evaluation for the CEM, PHX, CMUP, CMX and EMC lepton categories. Figure~\\ref{fig:pretag_fit} shows the result: the template composition with the fractions returned by the maximization likelihood fit. Table~\\ref{tab:wjet_pretag_qcd_fit} summarizes $F^{Wjets}_{Pretag}$ and $F^{nonW}_{Pretag}$ extracted from the fit. \n\nThe pretag non-$W$ fraction, $F^{nonW}_{Pretag}$, is not used any more in the background estimate but the total non-$W$ pretag normalization, $N^{nonW}_{Pretag}$, is needed to cross check the kinematic of the pretag control sample. We derive it with the following equation:\n\\begin{equation}\n N^{nonW}_{Pretag} = F^{nonW}_{Pretag} \\cdot N^{Data}_{Pretag} + \\sigma^{EWK}\\textrm{;}\n\\end{equation}\nwhere $\\sigma^{EWK}$ is the change in the normalization of the $EWK$ backgrounds as returned by the Likelihood fit.\n\n\n\\begin{table}\\begin{center}\n\\begin{tabular}{cccccc}\n\\toprule\n Lepton & CEM & PHX & CMUP & CMX & EMC \\\\\\midrule\n$F^{nonW}_{Pretag}$ & $ 8.5 \\pm 0.2$\\% & $ 13.6\\pm 0.1 $\\% & $ 1.9\\pm 0.2 $\\% & $2.2 \\pm 0.2 $\\% & $5.9\\pm 0.2$\\%\\\\\n$F^{Wjets}_{Pretag}$ & $ 84.0 \\pm 0.4$\\% & $ 82.5 \\pm 0.6$\\% & $86.7 \\pm 0.7 $\\% & $ 87.9 \\pm 0.9$\\% & $ 78.7\\pm 0.9 $\\%\\\\\n\\bottomrule\n\\end{tabular}\n\\caption[$W +$Jets and Non-$W$ Pretag Sample Compositions]{$W+$ jets and non-$W$ fractional compositions of the pretag selection sample estimated by a maximum likelihood fit on the SVM output distributions reported in Figure~\\ref{fig:pretag_fit}. The statistical error of the fit is reported.}\\label{tab:wjet_pretag_qcd_fit}\n\\end{center}\\end{table}\n\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_PretagSVTSVT_CEM}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_PretagSVTSVT_PHX}\\\\\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_PretagSVTSVT_CMUP}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_PretagSVTSVT_CMX}\\\\\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_PretagSVTSVT_ISOTRK}\n\\caption[$W +$ Jets and Non-$W$ Pretag Sample Estimate]{Background templates composition with proportions returned by the maximization likelihood fit on the SVM pretag distribution for the different lepton categories: CEM (top left), PHX (top right), CMUP (center left), CMX (center right), EMC (bottom). An arrow indicates the selection cut value in each lepton category.}\\label{fig:pretag_fit}\n\\end{center}\n\\end{figure}\n\n\\subsection{$W +$ Jets Scale Uncertainty}\nThe simulation of the $W+$~jets sample should take into account one more effect, the uncertainty on the factorization and renormalization scale, $Q^2$, of the ALPGEN event generator. \n\nThe $Q^2$ of the generator can be seen as the momentum scale of the hard interaction. Technically this value is used in two contexts:\n\\begin{itemize}\n\\item in the evaluation of the PDFs of the hard interaction, the {\\em factorization scale};\n\\item in the calculation of the ME perturbative expansion of a given QCD process, the {\\em renormalization scale};\n\\end{itemize}\nA reasonable choice for it, in the $W+n$~partons MC, is given by the following equation: \n\\begin{equation}\\label{eq:scale}\n Q^2 = M^2_W + \\sum^{partons} p_T^2 \n\\end{equation}\nwhere $M_W$ is the $W$ boson mass, the sum extends over all the partons and $p_T$ is their transverse energy.\n\nHowever the $Q^2$ is not a physical observable, as it is an artifact of the perturbative approximation needed to solve QCD problems, therefore an appropriate uncertainty should be taken into account. \nTable~\\ref{tab:q2samples} shows the characteristics of two new sets of $W+n$~partons MC samples used for systematic variation and generated with the $Q^2$ parameter doubled and halved.\n\\begin{table}\\begin{center}\n\\begin{tabular}{lcccc}\n\\toprule\n&\\multicolumn{2}{c}{\\bf $\\mathbf{Q^2=2.0}$ (Up)}& \\multicolumn{2}{c}{\\bf $\\mathbf{Q^2=0.5}$ Down}\\\\\nSample& $\\sigma$ (pb)& Initial Events & $\\sigma$ (pb)& Initial Events\\\\\n\\midrule\n$W+0p$ & 1767 & 3.5 M & 1912 & 3.5 M\\\\\n$W+1p$ & 182.9 & 3.5 M & 303 & 3.5 M\\\\\n$W+2p$ & 24.6 & 0.7 M & 57.6 & 0.7 M\\\\\n$W+3p$ & 3.36 & 0.7 M & 10.8 & 0.7 M\\\\\n$W+4p$ & 0.54 & 0.7 M & 2.24 & 0.7 M\\\\\n$W+b\\bar{b}+0p$& 2.30 & 0.7 M & 4.14 & 0.7 M\\\\\n$W+b\\bar{b}+1p$& 0.550 & 0.7 M & 1.64 & 0.7 M\\\\\n$W+b\\bar{b}+2p$& 0.152 & 0.7 M & 0.615 & 0.7 M\\\\\n$W+c\\bar{c}+0p$& 3.89 & 1.4 M & 6.88 & 1.4 M\\\\\n$W+c\\bar{c}+1p$& 1.08 & 1.4 M & 3.18 & 1.4 M\\\\\n$W+c\\bar{c}+2p$& 0.323 & 1.4 M & 1.30 & 1.4 M\\\\\n$W+c+0p$ & 13.8 & 1.4 M & 23.3 & 1.4 M\\\\\n$W+c+1p$ & 2.30 & 1.4 M & 5.72 & 1.4 M\\\\\n$W+c+2p$ & 0.294 & 1.4 M & 1.03 & 1.4 M\\\\\n$W+c+3p$ & 0.042 & 1.4 M & 0.189 & 1.4 M\\\\\n\\bottomrule\n\\end{tabular}\n\\caption[ALPGEN $W + $Partons MC Samples for $Q^2$ Systematics]{\\texttt{ALPGEN} LO plus \\texttt{PYTHIA} (PS) MC samples used to estimate the $Q^2$ systematic variation of the $W+HF$ and $W+LF$ backgrounds. The $Q^2$ is halved and doubled with respect to the default value given by Equation~\\ref{eq:scale}. Different MC samples are used for the three lepton flavors ($e, \\mu, \\tau$) of the $W$ decay.}\\label{tab:q2samples}\n\\end{center}\\end{table}\n\nThe $Q^2$ variation affects both the kinematic and the parton composition of the $W+$~jets sample, therefore, for each $Q^2$ systematic variation, the background estimate was redone as described Section~\\ref{sec:wjets_pretag}. The results are different shapes and rates for the $W+$~jets components. Figure~\\ref{fig:wjets_mjj_q2} shows the effect on the single-tag $M_{Inv}(jet1,jet2)$ distributions.\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/mjj_wlf_q2sys}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/mjj_wc_q2sys}\\\\\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/mjj_wbb_q2sys}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/mjj_wcc_q2sys}\n\\caption[$M_{Inv}(jet1,jet2)$ Distribution and $Q^2$ Variations for $W + ~LF\/c\/c\\bar{c}\/b\\bar{b}$ Samples]{$M_{Inv}(jet1,jet2)$ distribution for single \\texttt{SecVtx} tagged events, after combination of all the lepton categories: the default estimate of the four $W+$~jets samples is marked by the continuous line while the $Q^2$ systematic variations are marked by the dashed lines. Top left: $W+LF$. Top right: $W+c$. Bottom left: $W+c\\bar{c}$. Bottom right: $W+b\\bar{b}$.}\\label{fig:wjets_mjj_q2}\n\\end{center}\n\\end{figure}\n\n\n\n\\section{$W$ plus Light Flavors}\\label{sec:mistag}\n\nThe $W+$~jets events originating from a $LF$ quark can produce a secondary vertex for several reasons. Long living $LF$ hadrons produce a small amount of real \\verb'SecVtx' tags while false \\verb'SecVtx' tags are due to track reconstuction errors or interaction with the detector material and the beam pipe.\n\nThe $LF$ pretag fraction ($f_j^{LF}$) is derived from the $W+n$~partons MC composition, in the same way as the $HF$ fraction. However after requiring a $b$-tag, only a very small fraction of $LF$ jets remains in the sample. As the effects that generate mistags are not adequately simulated, appropriate parametrization is obtained studying the {\\em mistag} behaviour of the \\verb'SecVtx' algorithm. \n\nAs explained in Section~\\ref{sec:btag_sf}, a multi-jet control sample is used to parametrize a per-jet mistag probability, $p^j_{Mistag}$, function of six variables specific of the event and the jet. The $p^j_{Mistag}$ evaluation can be applied to any data sample or to MC. This gives the possibility to divide the normalization and the shape evaluation issues:\n\\begin{itemize}\n\\item the normalization is derived directly from the selected data sample. For each event we calculate an {\\em event-mistag} estimate, $w_{Mistag}^{ev}$, by using Equations~\\ref{eq:1tag} and~\\ref{eq:2tag} with the $p^j_{Mistag}$ of the jets as inputs. The sum of all the event-mistag gives a raw normalization of the total $LF$ contribution:\n \\begin{equation}\n N^{rawLF}=\\sum_{ev} w_{Mistag}^{ev}\\textrm{,}\n \\end{equation}\nthat needs to be corrected for the contribution of the other backgrounds. The $W+LF$ only part is:\n\\begin{equation}\n N^{W+LF} = N^{rawLF} \\frac{N^{Wjets} - N^{W+HF}- \\sum\\limits_{EWK}{N^{EWK}}}{N^{Wjets}+ \\sum\\limits_{EWK}N^{EWK}}\\textrm{.}\n\\end{equation}\n\n\\item Shape is derived from the $W+LF$ composition of ALPGEN MC (no $HF$ matched jets are used) where each simulated event is weighted by its own $w_{Mistag}$ to reproduce the kinematic beaviour of the mistagged jets. Figure~\\ref{fig:wlf_weight} shows the normalized ratios between the single-tagged $W+LF$ weighted MC and the original pretag $W+LF$ distribution: $LF$ jets have a higher mistag probability in the central $\\eta$ region of the detector (where also \\texttt{SecVtx} $b$-tagging efficiency is higher) and for large $E_T$. The effect on the $W+LF$ $M_{Inv}(jet1, jet2)$ spectrum after the $b$-tagging requirement is also shown.\n\n\\end{itemize}\n\nA rate uncertainty of $11$\\% and $21$\\% is included in the single and double-tagged mistag estimates respectively. It is derived from the mistag matrix parametrization and takes into account the statistical uncertainties and correlations between jets which fall in the same jet-bin of the mistag matrix. \n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/wlf_eta}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/wlf_et}\\\\\n\\includegraphics[width=0.7\\textwidth]{Bkg\/Figs\/wlf_mjj}\n\\caption[Normalized Ratios of $W+LF$ Pretag and Tag Distributions]{Normalized ratios between the single-tagged $W+LF$ weighted MC and the original pretag $W+LF$ MC distribution: jet 1 $\\eta$ distribution (top left), jet 1 $E_T$ distribution (top right) and $M_{Inv}(jet1, jet2)$ (bottom).}\\label{fig:wlf_weight}\n\\end{center}\n\\end{figure}\n\n\n\n\n\\section{Multi-jet Background}\\label{sec:qcd}\n\nAnother source of background comes from fake-$W$ events, where a QCD multi-jet event fakes the charged lepton and produces enough \\mbox{${\\not}{E_T}$} to pass the event selection (\\mbox{${\\not}{E_T}$} $>15$~GeV). \n\nThe probability to fake the lepton identification both at trigger and at offline level selection is very small but, at hadron colliders, multi-jet events are produced at such a high rate that this background becomes important. The selection of a fake lepton can happen for a mixture of physics and detector effects: for example a jet with large EM fraction and low track multiplicity can easily fake an electron or a high $p_T$ hadron can reach the muon chambers leaving little energy in not well instrumented sections of the calorimeter. A simulation of these effects would need an extremely high statistics, a perfect simulation of the detector and an accurate QCD prediction at high orders. This is not feasible therefore we rely on data to model the multi-jet sample (both shape and normalization). \n\nA fake-$W$ sample is obtained by reversing one or more $W\\to \\ell \\nu$ identification criteria so that the final sample will be enriched in multi-jet events and orthogonal to the signal selection. However the correlation between the reversed identification variables and the related kinematic quantities are lost in the process. The fake-$W$ sample may present discrepancies with respect to the truly multi-jet events selected in the signal region. \n\nThe production mechanism is strictly related to the lepton identification algorithm, therefore we employ three different models:\n\\begin{description}\n\\item[Not isolated muons:] the $IsoRel <0.1$ cut is one of the most important $W$ identification requirements therefore we can employ real muons but in the sideband with $IsoRel>0.2$ to select fake-$W$ candidates for CMX, CMUP and EMC categories\\footnote{EMC is composed by 7 lepton identification algorithms but the non-isolated EMC category is build only by four of them: CMU, CMP, CMXNT and BMU. The multi-jet sample properties are reproduced well enough also with this simpler model.}.\n For a better simulation, the jet associated with the non-isolated muon is removed from the jet-multiplicity count and its energy (without the muon) accounted for the \\mbox{${\\not}{E_T}$} correction (Section~\\ref{sec:metObj}).\n\\item[Fake-CEM:] CEM electron case is different because the isolation cut is correlated with the EM cluster energy, used both at trigger level selection and in the calorimeter \\mbox{${\\not}{E_T}$} calculation. A fake-$W$ model that maintains $IsoRel<0.1$ is obtained reversing at least two out of five shower identification requirements: Table~\\ref{tab:antiEle} shows the reversed cuts. As in the muon case, jet multiplicity should be corrected for the jet associated with the lepton, however it is not straightforward to decide if any correction needs to be applied to the \\mbox{${\\not}{E_T}$} calculation. \n\nWe studied the behaviour of the fake model on data, after full selection, but with no SVM cut and for \\mbox{\\mbox{${\\not}{E_T}$} $<15$~GeV}. The region is supposed to be dominated by multi-jet events (except for a small $Z$ contamination) therefore we compared data to the fake-$W$ model. Figure~\\ref{fig:chi2} shows how the reduced $\\chi^2$ test changes when comparing $M^W_T$ shapes for data and fake-$W$ model with different \\mbox{${\\not}{E_T}$} corrections. The best $\\chi^2$ value is obtained when the \\mbox{${\\not}{E_T}$} is corrected as if the jet corresponding to the anti-CEM is present in the event but with an energy of $0.45\\cdot E_T^{jet,raw}$. We apply this correction.\n\n\\item[Fake-PHX:] also the fake-$W$ PHX model is obtained with an anti-PHX selection where 2 out of 5 shower identification cuts are reversed (listed in Table~\\ref{tab:antiEle}). In this case no special prescription to the \\mbox{${\\not}{E_T}$} correction was found to drastically improve the kinematic model. Although some properties of the multi-jet sample are not perfectly reproduced, the high rejection power of the SVM discriminant reduces this to a minor problem.\n\\end{description}\nThe accurate description of the multi-jet background and its rejection by using the SVM discriminant played a major role in this analysis and in the latest $WH$ search results~\\cite{wh75_note10596,wh94_note10796}.\n\n\\begin{table}\\begin{center}\n\\begin{tabular}{c|c}\n\\toprule\n\\multicolumn{2}{c}{\\bf Anti-Electron: $\\ge 2$ Failed Cuts }\\\\\nCEM & PHX \\\\\\midrule\n$E^{Had}\/E^{EM}<0.055+0.0045 E^{EM}$& $E^{Had}\/E^{EM}<0.05$ \\\\ \n$L_{shr}<0.2$ & PEM$3\\times 3$Fit $\\ne 0$ \\\\ \n$\\chi^2_{\\mathrm{{CES strip}}}<10$ & $\\chi^2_{PEM3\\times 3} <10$ \\\\ \n$\\Delta z(CES,trk)<3$~cm & $PES5\\mathrm{by}9U>0.65$ \\\\ \n$-3.00.65$ \\\\ \n\\bottomrule\n\\end{tabular}\n\\caption[CEM and PHX Multi-jet Model Selection Cuts]{Multi-jet models for CEM and PHX electrons should fail at least 2 of the shower identification cuts listed here. The models built in this way are named {\\em anti-CEM} and {\\em anti-PHX} leptons.}\\label{tab:antiEle}\n\\end{center}\\end{table}\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.7\\textwidth]{Bkg\/Figs\/cem_scan100}\n\\caption[Reduced $\\chi^2$ Used to Tune \\mbox{${\\not}{E_T}$} Correction on Fake-$W$ Model]{Reduced $\\chi^2$ obtained comparing the normalized $M^W_T$ shapes of data and anti-CEM fake-$W$ model in the multi-jet enriched region \\mbox{${\\not}{E_T}$} $<15$~GeV. The test changes as we correct the \\mbox{${\\not}{E_T}$} for fractions of the energy associated with the jet corresponding to the anti-CEM.}\\label{fig:chi2}\n\\end{center}\n\\end{figure}\n\nThe fake-$W$ models can now be used in the background evaluation. Section~\\ref{sec:wjets_pretag} already gives the formula for the pretag $N^{nonW}$ estimate, however this result is not used in this analysis except for control purposes and to plot the background composition in the pretag region. On the other hand, the tagged region normalization of the fake-$W$ sample is the only missing piece in the full background evaluation. To estimate it we use again the multi-jet enriched sideband of the SVM output distribution.\n\nWe perform another maximum likelihood fit similar to the one described in Section~\\ref{sec:wjets_pretag} with two differences. The first is a modification of the template: we require events with at least $1$ ($2$) {\\em taggable} jets when evaluating the single (double) tag fake-$W$ templates, this allows an approximate simulation of the $b$-tag requirement retaining most of the statistics of the samples. The second difference is in the fit strategy: all the real-$W$ samples are added into one template whose normalization is left free to float in the fit together with $N^{nonW}_{k-Tag}$. This procedure and the use of the sidebands make the fake-$W$ determination as much uncorrelated as possible from the other backgrounds and from a possible unknown signal. Figures~\\ref{fig:nonw_tag1} and~\\ref{fig:nonw_tag2} show the results of the $N^{nonW}_{k-Tag}$ fit on the single and double tagged SVM distributions for the CEM, PHX, CMUP, CMX and EMC lepton categories.\nA conservative rate uncertainty of 40\\% is applied uniformly on all the lepton categories to account for low statistics in some of the fits and for the extrapolation from the sideband region. Fake-$W$ contamination is summarized in Table~\\ref{tab:wjet_tag_qcd_fit}, confirming the excellent rejection power of the SVM discriminant.\n\\begin{table}\\begin{center}\n\\begin{tabular}{cccccc}\n\\toprule\n Lepton & CEM & PHX & CMUP & CMX & EMC \\\\\\midrule\n$F^{nonW}_{1-Tag}$ & $ 9.7 \\pm 0.5$\\% & $ 19.1 \\pm 0.7 $\\% & $ 3.7 \\pm 0.7 $\\% & $ 3.3\\pm 1.0 $\\% & $ 6.2 \\pm 0.5$\\%\\\\\n $F^{nonW}_{2-Tag}$ & $ 3.2 \\pm 1.8$\\% & $ 10.8 \\pm 2.4$\\% & $ 4.7 \\pm 3.6$\\% & $ 0.0 \\pm 2.2$\\% & $ 0.0\\pm 0.3$\\%\\\\\n\\bottomrule\n\\end{tabular}\n\\caption[Non-$W$ Background Contamination in Single and Double Tag Samples]{Estimate of the non-$W$ contamination in the single and double-tagged selection samples. Only the statistical error of the fit is reported, a systematic uncertainty of $40$\\% is used in the final estimate.}\\label{tab:wjet_tag_qcd_fit}\n\\end{center}\\end{table}\n\n\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVT_CEM}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVT_PHX}\\\\\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVT_CMUP}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVT_CMX}\\\\\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVT_ISOTRK}\n\\caption[Single Tag SVM Distribution for Non-$W$ Background Estimate]{Result of the maximum likelihood fit performed on the single-tagged SVM distribution to estimate the non-$W$ background normalization: $N^{nonW}_{1-Tag} = F^{nonW}_{1-Tag}\\cdot N^{Data}_{k-Tag}$. Results are reported for the different lepton categories: CEM (top left), PHX (top right), CMUP (center left), CMX (center right), EMC (bottom). An arrow indicates the used selection cut value in each lepton category.}\\label{fig:nonw_tag1}\n\\end{center}\n\\end{figure}\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVTSVT_CEM}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVTSVT_PHX}\\\\\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVTSVT_CMUP}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVTSVT_CMX}\\\\\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/NonWFit_2JET_SVTSVT_ISOTRK}\n\\caption[Double Tag SVM Distribution for Non-$W$ Background Estimate]{Result of the maximum likelihood fit performed on the double-tagged SVM distribution to estimate the non-$W$ background normalization: $N^{nonW}_{1-Tag} = F^{nonW}_{1-Tag}\\cdot N^{Data}_{k-Tag}$. Results are reported for the different lepton categories: CEM (top left), PHX (top right), CMUP (center left), CMX (center right), EMC (bottom). An arrow indicates the used selection cut value in each lepton category.}\\label{fig:nonw_tag2}\n\\end{center}\n\\end{figure}\n\n\\section{Correction of Luminosity Effects}\n\nThe effect of instantaneous luminosity should be corrected in the available MC samples before completing of the background estimate. \n\nThe instantaneous luminosity of the collisions influences several parameters like the occupancy of the detector, the number of vertices in the interactions and so on. Most of these effects are parametrized within the CDF simulation (see Section~\\ref{sec:mcGen}) and they are simulated in a run-dependent way (i.e. each MC section corresponds to a period of data taking with appropriate tuning). Unluckily the available MC samples reproduce only a first section (up to 2008) of the CDF dataset while the highest luminosities (easily above $\\mathscr{L} = 2\\times 10^{32}$~cm$^{-2}$s$^{-1}$) were delivered by the Tevatron in the latest years. \n\nAn approximate solution to the problem is the {\\em reweighting} of the MCs as a function of the observed number of good quality primary vertices\\footnote{Quality$\\ge 12$, see Table~\\ref{tab:pVtx}.} (a variable highly correlated with the instantaneous luminosity). The reweighting function is evaluated comparing, in the pretag control region, the number of vertices estimated in MC against the one observed in data. Figure~\\ref{fig:nvtx_rew} shows the distribution of the number of primary vertices withouth any correction and after correcting the background estimate with the reweighting function that we apply to the final evaluation.\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_nQ12Vtx_noW}\n\\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_nQ12Vtx}\n\\caption[Number of \\mbox{{\\em Quality} $\\ge 12$} Before and After Reweighting]{Distribution (for all the lepton categories) of the number of \\mbox{{\\em Quality} $\\ge 12$} primary vertices for the events selected in the pretag control region before (left) and after (right) the application of a reweighting function to the MC events.}\\label{fig:nvtx_rew}\n\\end{center}\n\\end{figure}\n\n\n\n\n\n\n\\section{Final Background Evaluation}\n\nAs explained in the previous part of this Chapter, we derive the shape and normalization contributions of each background in an independent way. The composition of them should reproduce, within uncertainties, the observed number of selected events and the kinematic distributions of data in control and signal regions. If the background estimate is proven to be solid, then it is possible to investigate the presence of a signal with adequate tools. \n\nWe present here the final background estimate. Statistical analysis and the measured properties of the $diboson\\to \\ell \\nu + HF$ signal will be described in the next Chapter.\n\nThe first important validation criteria is the agreement of the kinematic distributions in the pretag control region. As the normalization of the region is constrained by data the most important effect comes from the evaluation of the shapes. Figures from~\\ref{fig:pret_jetEt} to~\\ref{fig:pret_mjj} show several kinematic variables after the composition of all the background for all the lepton categories. For each variable, we estimated a central pretag background value and four systematic variations: \n\\begin{itemize}\n\\item jet energy scale plus and minus one sigma (JES $+1\\sigma$, JES $-1\\sigma$);\n\\item $W+n$~partons renormalization scale variation: $Q^2$ is doubled ($Q^2_{Up}= 2.0 Q^2 $) and halved ($Q^2_{Down}= 0.5 Q^2$).\n\\end{itemize}\nThe positive (negative) systematic variations are added in quadrature and the result is shown together with the central background evaluation. The agreement is excellent within uncertainties. Few distributions that show mis-modeling, like the $\\eta$ of the second jet shown in Figure~\\ref{fig:pret_jetEta}, become irrelevant after the tagging requirement that enforces a more central selection of the jets.\n\n \\begin{figure}\n \\begin{center} \n \\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_J1Et}\n \\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_J2Et}\n \\end{center}\n \\caption[Jets $E_T$ Distributions in Pretag Control Region]{Jets corrected $E_T$ distributions in the pretag control region, for all the lepton categories combined: jet 1 $E_T$ (left) and jet 2 $E_T$ (right) for all the lepton categories combined, in the pretag control region.}\\label{fig:pret_jetEt}\n \\end{figure}\n\n \\begin{figure}\n \\begin{center} \n \\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_J1Eta}\n \\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_J2Eta}\n \\end{center}\n \\caption[Jets $\\eta$ Distributions in Pretag Control Region]{Jets $\\eta$ distributions in the pretag control region, for all the lepton categories combined: jet 1 $\\eta$ (left) and jet 2 $\\eta$ (right).}\\label{fig:pret_jetEta}\n \\end{figure}\n \n \\begin{figure}\n \\begin{center} \n \\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_LepPt}\n \\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_LepEta}\n \\end{center}\n \\caption[Lepton $p_T$ and $\\eta$ Distributions in Pretag Control Region]{Two lepton related kinematic distributions in the pretag control region, for all the lepton categories combined: $p_T$ (left) and $\\eta$ (right).}\\label{fig:pret_lepKin}\n \\end{figure}\n\n \\begin{figure}\n \\begin{center}\n \\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_met}\n \\includegraphics[width=0.495\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_dPhiMetJet1}\n \\end{center}\n \\caption[$\\nu$ Kinematic Distributions in Pretag Control Region]{Two $\\nu$ related kinematic distributions in the pretag control region, for all the lepton categories combined: \\mbox{${\\not}{E_T}$} (left) and $\\Delta\\phi(jet1, $\\mbox{${\\not}{E_T}$} $)$ (right).}\\label{fig:pret_metKin}\n \\end{figure}\n\n \\begin{figure}\n \\begin{center} \n \\includegraphics[width=0.49\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_wPt}\n \\includegraphics[width=0.49\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_WMt}\n \\caption[$W$ Kinematic Distributions in Pretag Control Region]{Two $W$ related kinematic distributions in the pretag control region, for all the lepton categories combined: $p_T^W$ (left) and $M_T^W$ (right).}\\label{fig:pret_wkin}\n \\end{center}\n \\end{figure}\n\n \\begin{figure}\n \\begin{center}\n \\includegraphics[width=0.49\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_dRLepJ1}\n \\includegraphics[width=0.49\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_dRJ1J2} \n \\end{center}\n \\caption[Angular Separation Distributions in Pretag Control Region]{Two angular separation distributions in the pretag control region, for all the lepton categories combined: $\\Delta R(lep,jet1)$ (left), $\\Delta R(jet1,jet2)$, (right).}\\label{fig:pret_dphi}\n \\end{figure}\n\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=0.98\\textwidth]{Bkg\/Figs\/WZ_2JET_PretagSVT_CEMCMUPCMXISOTRKPHX_Mj1j2}\n\\caption[$M_{Inv}(jet1, jet2)$ Distribution in Pretag Control Region]{Invariant mass distribution of the selected di-jet pair in the pretag control region, for all the lepton categories combined.}\\label{fig:pret_mjj}\n\\end{center}\n\\end{figure}\n\nThe other validation method is the actual comparison of the expected backgrounds in the signal region against the observed selection. Tables~\\ref{tab:1tag_final} and~\\ref{tab:2tag_final} report the complete background composition in each of the separate lepton categories for the single and double \\texttt{SecVtx}-tagged selections, neither JES nor $Q^2$ variations are reported in the tables. All the estimates are within the uncertainties, that are, however, rather large because of the $K^{HF}$ factor uncertainties listed in Table~\\ref{tab:k-fac}. \n\n\\begin{sidewaystable}\n \\begin{center}\n \\begin{tabular}{lccccc}\n \\toprule\n \\multicolumn{6}{c}{\\bf Single-tag Event Selection}\\\\\n Lepton ID & CEM & PHX & CMUP & CMX & EMC \\\\ \n\\midrule \n $Z+$jets & 55.53 $\\pm$ 4.73 & 7.76 $\\pm$ 0.67 & 65.3 $\\pm$ 5.73 & 37.14 $\\pm$ 3.28 & 104.18 $\\pm$ 10.8 \\\\ \n $t\\bar{t}$ & 237.55 $\\pm$ 23.3 & 46.93 $\\pm$ 4.59 & 139.68 $\\pm$ 13.8 & 62.22 $\\pm$ 6.14 & 228.58 $\\pm$ 25.6 \\\\ \n Single-top $s$ & 64.23 $\\pm$ 5.88 & 11.08 $\\pm$ 1.02 & 36.42 $\\pm$ 3.35 & 16.11 $\\pm$ 1.49 & 51.68 $\\pm$ 5.48 \\\\ \n Single-top $t$ & 84.95 $\\pm$ 9.99 & 16.86 $\\pm$ 1.98 & 47.75 $\\pm$ 5.64 & 22.5 $\\pm$ 2.66 & 66.22 $\\pm$ 8.56 \\\\ \n $WW$ & 84.35 $\\pm$ 11.8 & 25.05 $\\pm$ 3.54 & 43.7 $\\pm$ 6.17 & 23.68 $\\pm$ 3.35 & 53.11 $\\pm$ 8 \\\\\n $ZZ$ & 1.85 $\\pm$ 0.19 & 0.21 $\\pm$ 0.02 & 2.45 $\\pm$ 0.25 & 1.37 $\\pm$ 0.14 & 3.45 $\\pm$ 0.39 \\\\ \n $WZ$ & 29.2 $\\pm$ 2.95 & 12.32 $\\pm$ 1.22 & 16 $\\pm$ 1.66 & 9.21 $\\pm$ 0.94 & 20.54 $\\pm$ 2.39 \\\\\n $W+b\\bar{b}$ & 858.71 $\\pm$ 258 & 263.08 $\\pm$ 79.1 & 428.58 $\\pm$ 129 & 238.99 $\\pm$ 71.9 & 409.97 $\\pm$ 123 \\\\\n $W+c\\bar{c}$ & 441.77 $\\pm$ 134 & 145.08 $\\pm$ 44.1 & 212.98 $\\pm$ 64.8 & 120.26 $\\pm$ 36.6 & 214 $\\pm$ 65 \\\\\n $W+cj$ & 342.86 $\\pm$ 104 & 96.42 $\\pm$ 29.3 & 171.77 $\\pm$ 52.2 & 93.01 $\\pm$ 28.3 & 143.47 $\\pm$ 43.6 \\\\ \n $W+LF$ & 809.01 $\\pm$ 87.1 & 302.78 $\\pm$ 32.2 & 408.74 $\\pm$ 43.3 & 230.75 $\\pm$ 24.7 & 463.65 $\\pm$ 53.5 \\\\\n Non-$W$ & 302.69 $\\pm$ 121 & 205.06 $\\pm$ 82 & 58.75 $\\pm$ 23.5 & 27.96 $\\pm$ 11.2 & 106.35 $\\pm$ 42.5 \\\\\n\\midrule \n {\\bf Prediction} & {\\bf 3312.68 $\\pm$ 521} & {\\bf 1132.63 $\\pm$ 176 } & {\\bf 1632.11 $\\pm$ 253 } & {\\bf 883.21 $\\pm$ 140} & {\\bf 1865.2 $\\pm$ 248} \\\\\n {\\bf Observed } & {\\bf 3115} & {\\bf 1073} & {\\bf 1577 } & {\\bf 830} & {\\bf 1705} \\\\\n\\midrule \n {\\bf Dibosons} & {\\bf 115.39 $\\pm$ 13} & {\\bf 37.57 $\\pm$ 4.05} & {\\bf 62.15 $\\pm$ 6.87} & {\\bf 34.26 $\\pm$ 3.75} & {\\bf 77.09 $\\pm$ 9.35} \\\\\n\\bottomrule\n \\end{tabular} \n \\end{center}\n \\caption[Observed and Expected Events with One \\texttt{SecVtx} Tags]{Summary of observed and expected events with one secondary vertex tag (\\texttt{SecVtx}), in the $W+2$~jets sample, in \n $9.4$~fb$^{-1}$ of data in the different lepton categories. Statistical and systematics rate uncertainties are included in the table except for the contribution of the JES and the $Q^2$ variations. Diboson expected yield is added in the prediction.}\\label{tab:1tag_final}\n\\end{sidewaystable}\n\n\\begin{sidewaystable}\n \\begin{center}\n\n \\begin{tabular}{lccccc}\n \\toprule \n \\multicolumn{6}{c}{\\bf Double-tag Event Selection}\\\\\n Lepton ID & CEM & PHX & CMUP & CMX & EMC \\\\ \n \\midrule \n \n \n $Z+$jets & 1.43 $\\pm$ 0.13 & 0.27 $\\pm$ 0.02 & 2.84 $\\pm$ 0.27 & 1.41 $\\pm$ 0.13 & 4.6 $\\pm$ 0.5 \\\\\n $t\\bar{t}$ & 48.21 $\\pm$ 6.99 & 9.91 $\\pm$ 1.44 & 27.31 $\\pm$ 3.98 & 12.37 $\\pm$ 1.8 & 44.89 $\\pm$ 6.95 \\\\\n Single-top $s$ & 16.89 $\\pm$ 2.36 & 2.87 $\\pm$ 0.4 & 9.68 $\\pm$ 1.36 & 4.17 $\\pm$ 0.58 & 13.72 $\\pm$ 2.05 \\\\\n Single-top $t$ & 5.07 $\\pm$ 0.81 & 1.13 $\\pm$ 0.18 & 2.85 $\\pm$ 0.46 & 1.34 $\\pm$ 0.22 & 4.16 $\\pm$ 0.71 \\\\\n $WW$ & 0.72 $\\pm$ 0.19 & 0.18 $\\pm$ 0.05 & 0.35 $\\pm$ 0.09 & 0.2 $\\pm$ 0.05 & 0.49 $\\pm$ 0.13 \\\\\n $ZZ$ & 0.26 $\\pm$ 0.04 & 0.03 $\\pm$ 0 & 0.46 $\\pm$ 0.06 & 0.29 $\\pm$ 0.04 & 0.63 $\\pm$ 0.1 \\\\ \n $WZ$ & 5.28 $\\pm$ 0.75 & 2.6 $\\pm$ 0.37 & 2.52 $\\pm$ 0.36 & 1.67 $\\pm$ 0.24 & 3.52 $\\pm$ 0.54 \\\\\n $W+b\\bar{b}$ & 114.7 $\\pm$ 35.1 & 33.92 $\\pm$ 10.4 & 59.06 $\\pm$ 18.1 & 29.49 $\\pm$ 9.04 & 60.71 $\\pm$ 18.6 \\\\\n $W+c\\bar{c}$ & 6.68 $\\pm$ 2.1 & 2.16 $\\pm$ 0.68 & 3.41 $\\pm$ 1.08 & 1.63 $\\pm$ 0.51 & 4 $\\pm$ 1.25 \\\\ \n $W+cj$ & 5.18 $\\pm$ 1.63 & 1.43 $\\pm$ 0.45 & 2.75 $\\pm$ 0.87 & 1.26 $\\pm$ 0.4 & 2.69 $\\pm$ 0.84 \\\\ \n $W+LF$ & 4.53 $\\pm$ 0.94 & 1.7 $\\pm$ 0.36 & 2.35 $\\pm$ 0.48 & 1.28 $\\pm$ 0.27 & 2.98 $\\pm$ 0.66 \\\\ \n Non-$W$ & 5.58 $\\pm$ 2.23 & 6.79 $\\pm$ 2.72 & 4.31 $\\pm$ 1.73 & 0 $\\pm$ 0.5 & 0 $\\pm$ 0.5 \\\\ \n \\midrule \n {\\bf Prediction} & {\\bf 214.53 $\\pm$ 40.5} & {\\bf 62.99 $\\pm$ 12.1} & {\\bf 117.92 $\\pm$ 21.1} & {\\bf 55.11 $\\pm$ 10.4} & {\\bf 142.38 $\\pm$ 23.3} \\\\\n {\\bf Observed} & {\\bf 175} & {\\bf 62 } & {\\bf 92} & {\\bf 49 } & {\\bf 126} \\\\ \n \\midrule \n {\\bf Dibosons } & {\\bf 6.26 $\\pm$ 0.79} & {\\bf 2.8 $\\pm$ 0.38} & {\\bf 3.34 $\\pm$ 0.4} & {\\bf 2.15 $\\pm$ 0.26} & {\\bf 4.64 $\\pm$ 0.61 } \\\\\n \\bottomrule\n \\end{tabular} \n \\end{center}\n \\caption[Observed and Expected Events with Two \\texttt{SecVtx} Tags]{Summary of observed and expected events with two secondary vertex tags (\\texttt{SecVtx}), in the $W+2$~jets sample, in \n $9.4$~fb$^{-1}$ of data in the different lepton categories. Statistical and systematics rate uncertainties are included in the table except for the contribution of the JES and the $Q^2$ variations. Diboson expected yield is added in the prediction.}\\label{tab:2tag_final}\n\\end{sidewaystable}\n\nBackground is still large although the diboson signal is sizable. An analysis of the $diboson\\to \\ell\\nu+HF$ properties with a counting experiment is not feasible. Therefore, in the next Chapter, we exploit the separation power of two variables: \n\\begin{itemize}\n\\item the di-jet invariant mass, $M_{Inv}(jet1,jet2)$, improves the separation between the non-resonant $W+HF$ production and the diboson.\n\\item KIT-NN improves the $b$-jets $vs$ $c$-jets separation.\n\\end{itemize}\nThe shape analysis and the fitting procedure also constraint the large normalization uncertainties, thereby increasing the significance of the measurement.\n\n\n\n\n\\chapter{Conclusions}\\label{chap:conclusions}\n\nThe result of this thesis is the evidence, at $3.08\\sigma$, for the associate production of massive vector bosons ($W$ and $Z$) detected at the CDF II experiment in a final state with one lepton, \\mbox{${\\not}{E_T}$} and $HF$-tagged jets. \n\nSuch a result, obtained on the complete CDF II dataset ($9.4$~fb$^{-1}$ of data), was possible thanks to a simultaneous effort in several directions. The signal acceptance was extended both at online and offline selection level; the multi-jet background was strongly suppressed using a SVM-based multi-variate algorithm; second vertex $HF$-tagging was used in conjunction with a flavor-separator NN (KIT-NN). \n\nThe signals, both inclusive diboson production and $WW$ {\\em vs} $WZ\/ZZ$ separately, were extracted from the invariant mass distribution, $M_{Inv}(jet1,jet2)$, of single and double $HF$-tagged jet pairs. For single-tagged events the $b$ quark {\\em vs} $c$ quark discrimination was obtained by using a bi-dimensional distribution $M_{Inv}(jet1,jet2)$ {\\em vs} KIT-NN.\n\nWe measure the total diboson production cross section, fixing the $WW$ and $WZ\/ZZ$ relative contribution to the SM prediction. We obtain:\n\\begin{equation}\n \\sigma_{Diboson}^{Obs}=\\left(0.79\\pm 0.28\\right)\\times \\sigma_{Diboson}^{SM} = \\left(14.6\\pm 5.2\\right)\\textrm{~pb,}\n\\end{equation}\nwhere the errors include statistical and systematic uncertainty and $\\sigma^{SM}_{Diboson}$ is the SM predicted cross section.\n\nThen, after removing the constraint on the $WW$ and $WZ\/ZZ$ relative contribution, we leave them free to float independently. We perform a simultaneous measurement of $\\sigma_{WW}$ and $\\sigma_{WZ\/ZZ}$ where all correlations are included. We can also perform a separate measurement of each contribution one at the time by considering the other as a background. In this case we obtain:\n\\begin{equation}\n\\sigma^{Obs}_{WW}=\\left(0.45^{ +0.35}_{ -0.32}\\right)\\times\\sigma^{SM}_{WW} = \\left(5.1^{+4.0}_{-3.6}\\right)\\textrm{~pb}\n\\end{equation}\nand\n\\begin{equation}\n \\sigma^{Obs}_{WZ\/ZZ}=\\left(1.64^{+0.83}_{-0.78}\\right)\\times\\sigma^{SM}_{WZ\/ZZ} = \\left(11.6^{+5.9}_{-5.5}\\right)\\mathrm{~pb,} \n\\end{equation}\nwhere the errors include statistical and systematic uncertainties.\nThe significance of the measurements is $1.78\\sigma$ for $WW$ signal and $2.54\\sigma$ for $WZ\/ZZ$ signal.\n\nAll the results are consistent with the SM prediction and they confirm the CDF capability of identify a small signal in this challenging final state. In particular the previous version of this analysis, performed with a dataset of $7.5$~fb$^{-1}$ and reported in Appendix~\\ref{App:7.5}, was the first evidence of diboson production in $\\ell\\nu+HF$ final state at a hadron collider.\n\nBeyond the pure testing of SM predicted processes, several of the techniques developed for this thesis were also applied to the $WH$ search at CDF, with a relevant improvement of the sensitivity to this process.\n\n\\chapter{The CDF II Experiment}\\label{chap:detector}\nIn this Chapter the accelerator facility and the detector apparatus are described in their main features. \n\nThe analyzed dataset presented in this thesis corresponds to $9.4$~fb$^{-1}$ of data collected by the CDF II (Collider Detector at Fermilab for Run II) experiment along its entire operation time\\footnote{Most of the presented work was performed during the data taking period therefore some experimental features are presented as if the operations are not yet concluded.}, from February 2002 to September $30^{th}$, 2011.\nThe CDF multi-purpose detector was located at one of the two instrumented interaction points along the Tevatron accelerator ring where $p\\bar{p}$ beams collided at an energy of $\\sqrt{s}=1.96$~TeV. \n\n\\section{The Tevatron}\n\nThe Tevatron collider was a proton-antiproton storage ring and circular accelerator located at the Fermi National Accelerator Laboratory (FNAL or Fermilab), $50$~Km west from Chicago (Illinois, U.S.A.). With a center-of-mass energy of $\\sqrt{s}=1.96$~TeV, it was the world highest energy accelerator~\\cite{tevatronweb} before the beginning of the Large Hadron Collider (LHC) era\\footnote{The Tevatron collision energy record was exceed by the LHC on March $30^{th}$, 2010, when the first $pp$ collisions at $\\sqrt{s}=7$~TeV took place.} and the largest anti-matter source in the world. The decommissioning of the accelerator started at the end of 2011 with the last collision and the stop of the operations on September $30^{th}$, 2011.\n\nThe history of the Tevatron is marked by impressive technology achievements and physics results. For example, starting the operations in 1983, it was the fist super-conducting magnet accelerator ring, in 1995, the {\\em top} quark was discovery here~\\cite{CDF_top_obs, d0_top_obs} or in 2006, $B_s$ mixing was observed~\\cite{cdf_bsmix_obs}.\n\nOne of most striking achievement of the Tevatron was the whole process of the proton-antiproton production and acceleration that involved the simultaneous operation of a chain of accelerator machines. Figure~\\ref{accel_sec} shows a view of Tevatron complex and of its sections~\\cite{tevatronweb,rookie_book}, next paragraphs summarize their operation and performances.\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{Detector\/Figs\/accel_sec}\n\\caption[CDF Detector Schema]{Schematic view of Tevatron accelerator complex at Fermilab, different colors mark different accelerator sections.}\\label{accel_sec}\n\\end{center}\n\\end{figure}\n\n\\subsection{Proton and Antiproton Production}\n\nThe first stage, proton extraction and initial acceleration, takes place in the Pre-Accelerator (\\emph{PreAc}). Hot hydrogen gas molecules ($H_{2}$) are \nsplit by an intense local electrostatic field and charged with two electrons; $H^{-}$ ions are accelerated up to $750$~KeV by a Cockcroft-Walton accelerator every $66$~ms.\n\nPreAc ion source\\footnote{The ion sources are actually two, named H- and \nI-, and working alternatively.} constantly produces beams at $15$~Hz rate and \nsend them to the \\emph{Linac}: a linear accelerator that increases the ions \nenergy $750$~KeV to $400$~MeV. It is made of two sections: a low energy \ndrift tube and a high energy coupled cavity at the end.\n\nNext acceleration stage is performed by a circular accelerator (synchrotron)\nof $75$~m radius called \\emph{Booster}. The insertion of a thin carbon foil \nstrips off electrons from the $400$~MeV ions and a sweep, from $\\sim38$ to \n$\\sim53$~MHz in radio-frequency (RF), carries resulting protons to an energy \nof $8$~GeV. The use of negative ions permits injection of more particles from \nthe Linac, otherwise the magnetic field needed to catch the protons would also\nkick away protons already inside the Booster. Bunches are extracted when about\n$8\\cdot 10^{12}$ protons are collected.\n\nResulting bunches have the correct energy to be sent to the \\emph{Main Injector}, a larger synchrotron with radius of $\\approx 0.5$~Km): here conventional magnets and 18 accelerating cavities are employed to accelerate protons up to $120$~GeV or $150$~GeV, depending upon their use. The $120$~GeV protons are used to produce the antiprotons, while $150$~GeV protons are further accelerated into the Tevatron main ring. \n\nAntiproton production takes place in the \\emph{Antiproton Source}. This machine is composed by several parts (see Figure~\\ref{accel_sec}): first \nthere is a target station where the $120$~GeV protons, extracted from Main Injector, collide with a Nickel target and $8$~GeV $\\bar{p}$ are selected from all the resulting particles. Typically, $10\\div 20$ $\\bar{p}$ are collected for each $10^{6}$ protons on target. After production, antiprotons have a large spatial and momentum spread while acceleration into Main Injector requires narrow $8$~GeV packets. Therefore they are sent to the \\emph{Debuncher}: a triangular shape synchrotron, with a mean radius of $90$~m. The bunch signal is picked up and analyzed at one side of the ring and then it is corrected on the other side in a process named \\emph{stochastic cooling} and \\emph{bunch rotation}~\\cite{Marriner:2003mn_cooling}.\nThe final step of the $\\bar{p}$ production is the accumulation: the beam is sent to a smaller synchrotron (with a mean radius of $75$~m) inside Debuncher ring called \\emph{Accumulator}. Other cooling methods are applied here.\n\nFrom the Accumulator, $8$~GeV antiprotons can be transferred either to the Main\nInjector or to the Recycler ring. The latter is a $3.3$~Km long ring of permanent magnets, located in the Main Injector enclosure, which is used to gather antiprotons as a final storage\nbefore the injection into the Tevatron, thus allowing the Accumulator to operate at its optimal efficiency.\n\nDuring the last year of running, 2011, a store could start with up to $3.5\\times 10^{12}$ antiprotons, collected in $10\\div 20$ hours of production.\n\n \n\\subsection{Collision and Performance}\n\nLast acceleration stage takes place into the \\emph{Tevatron} Main Ring: with a radius of one kilometer this is the largest of the Fermilab accelerators, and, thanks to superconducting magnets, it can store and \naccelerate beams from an energy of $150$~GeV (Main Injector result) to $980$~GeV.\nTable~\\ref{tab:FermilabAcceleratorComplex} summarizes the acceleration characteristics of the different stages of the Fermilab $p\\bar{p}$ Accelerator Complex.\n\n\\begin{table}[h]\n\\begin{center}\n\\begin{small}\n\\begin{tabular}{cccccccc}\\toprule\nAcc. & H & $H^{-}$ & C-W & L & B & M & T\\\\\n\\midrule\nE &0.04 eV & 25 KeV &750 KeV & 400 MeV&8 GeV& 150 GeV&0.98 TeV\\\\\n$\\beta$ & $9.1\\cdot10^{-8}$ & $0.01$& $0.04$ & $0.71$ & $0.99$ & $1$ &$1$\\\\\n$\\gamma$ & $1$ & $1$ & $1$ & $1.43$ & $9.53$ & $161$ & $1067$\\\\\n\\bottomrule\n\\end{tabular}\n\\end{small}\n\\caption[Performance of Fermilab Accelerator Complex]\n{Performances of the Fermilab Accelerator Complex. The steps along the accelerator chain (Acc.), with the corresponding labelling, are: Cockwroft-Walton (C-W), Linac (L), Booster (B), Debuncher and Recycler, Main Injector (M), Tevatron (T). The energy reached at the end of the step is E, $\\beta=\\frac{v}{c}$ expresses\nthe speed of the particle as a fraction of the speed of light in vacuum and \\mbox{$\\gamma=\\frac{E}{pc}=\\frac{1}{\\sqrt{1-(\\frac{v}{c})^{2}}}$} is the relativistic factor.}\\label{tab:FermilabAcceleratorComplex}\n\\end{center}\n\\end{table}\n\n\nWhen beams production and acceleration is complete, a Tevatron \\emph{store} is started: 36 protons and 36 antiprotons bunches, containing respectively \n$\\sim 10^{13}$ and $\\sim 10^{12}$ particles, are injected into the Main Ring at location\\footnote{The Tevatron is divided into six sections (see \nFigure~\\ref{accel_sec}) and each junction zone, named form A to F, has a different function: most important areas are B0, D0 and F0, the first two are experimental areas where CDF and DO~detectors are placed, while F0 contains RF cavity for beam acceleration and switch areas to connect Main Injector and the Tevatron.} F0 to be collided. The $0.5$~mm thin proton and antiproton bunches share the same beam pipe, magnets and vacuum system and they follows two non intersecting orbits kept $5$~mm away from each other. Beam control is obtained through nearly 1000 superconducting magnets $6$~m long, cooled to $4.3$~K and capable of $4.2$~T fields. \n\\\\\n\nBeside energy, the other fundamental parameter of an accelerator is the \ninstantaneous luminosity ($\\mathscr{L}$), as the rate of a physical process\nwith cross section $sigma$ is:\n\\begin{equation}\\label{ist_ev}\n\\frac{dN}{dt}[\\textrm{events s}^{-1}]=\\mathscr{L}[\\textrm{cm}^{-2}\\textrm{s}^{-1}]\\sigma[\\textrm{cm}^{2}]\\mathrm{.}\n\\end{equation}\nHigh energy permits an insight to incredibly small scale physics but only very \nhigh instantaneous luminosity and very large integrated (in time) luminosity\nallow to see rare events. Figure~\\ref{processes} shows the production\ncross section of different physical processes\\footnote{Due to the tiny cross sections we deal with, through most of this work we will be using {\\it picobarns} (pb) where $1$~pb$= 10^{-36}\\mathrm{~cm}^{2}$.}. \n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.75\\textwidth]{Detector\/Figs\/processes_new}\n\\caption[Production Cross Section of Physical Processes at CDF]{Predicted production cross section of physical processes at CDF along with their measured values. Higgs cross section varies with the Higgs mass.}\\label{processes}\n\\end{center}\n\\end{figure}\n\nThe instantaneous luminosity of an accelerator, usually measured in cm$^{-2}$s$^{-1}$, is given by:\n\\begin{equation}\\label{ist_lum}\n\\mathscr{L}=\\frac{N_{p}N_{\\bar{p}}B f}{2\\pi(\\sigma^2_{p}+\\sigma^2_{\\bar{p}})} F\\Big(\\frac{\\sigma_l}{\\beta^\\ast}\\Big),\n\\end{equation}\nwhere $N_{p}$ ($N_{\\bar{p}}$) are the number of protons (antiprotons) per bunch, $B$ is the number of bunches inside accelerator, $f$ is the bunch \ncrossing frequency, $\\sigma_{p(\\bar{p})}$ is the r.m.s. of the proton (antiproton) beam at the interaction point and $F$ is a beam shape form factor depending on the ratio between the the longitudinal r.m.s. of the bunch, $\\sigma_l$, and the beta function, $\\beta^\\ast$, a measure of the beam extension in the $x,y$ phase space.\nSeveral of these parameters are related to the accelerator structure so they are (almost) fixed inside the Tevatron:\n36 $p\\bar{p}$ bunches crossed with frequency of $396$~ns and, at interaction points, quadrupole magnets focus beams in $\\approx 30$~$\\mu$m spots with $\\beta^\\ast\\simeq 35$~cm and $\\sigma_l \\simeq 60$~cm.\n\nHigher luminosities were achieved in Run II thanks to the increased antiproton stack rate. Figures~\\ref{initLum} and~\\ref{lumi} show the delivered luminosity and the constant progress in the performances of the machine with an instantaneous luminosity record of $4.3107\\times 10^{32}$~cm$^{-2}$s$^{-1}$ on May, 3$^{rd}$ 2011~\\cite{acc_records}. \n\nThe total integrated luminosity produced by the Tevatron is more than $12$~fb$^{-1}$ and CDF wrote on tape, on average, about $85\\%$ of it, with small inefficiencies due\nto detector calibration, stores not used to collect data for physics or dead time during the start-up of the data taking. This analysis uses all the data collected by CDF in Run II, corresponding, after data quality requirements, to a integrated luminosity of about $9.4$~fb$^{-1}$.\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{Detector\/Figs\/initLum}\n\\caption[Initial Instantaneous Luminosity Delivered by the Tevatron to CDF]{Initial instantaneous luminosity delivered by the Tevatron accelerator to the CDF II detector. Performances grown by two order of magnitude from the beginning of the operation (2002) to the final collisions (2011)~\\cite{lum}.}\\label{initLum}\n\\end{center}\n\\end{figure}\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{Detector\/Figs\/lum}\n\\caption[Integrated Luminosity Delivered by the Tevatron for each Run]{Integrated luminosity delivered by the Tevatron for each physics run (blue) and averaged on a week by week basis (green bars)~\\cite{lum}.}\\label{lumi}\n\\end{center}\n\\end{figure}\n\n\n\n\\section{The CDF Detector}\n\nCDF II is a multi-purpose solenoidal detector situated at the B0 interaction point along the Tevatron main accelerator ring. Thanks to accurate charged \nparticle tracking, fast projective calorimetry and fine grained muon detection, the CDF II detector can measure energy, momentum and charge of most particles resulting from $\\sqrt{s}=1.96$~TeV $p\\bar{p}$ collisions. \n\nThe first original design goes back to 1981 but CDF underwent many upgrades \nduring the past twenty years. The last and most extensive one began in \n1996 and ended in 2001 when Tevatron \\emph{Run II} started. At present the \nCDF II experiment is operated by an international collaboration\nthat embraces more than 60 institutions from 13 different countries, for a total\nof about 600 researchers.\n\n\n\\subsection{Overview and Coordinate system}\n\nCDF is composed by many parts (sub-detectors) for a total of about \n5000 tons of metal and electronics, a length of $\\sim 16$~m and a diameter of \n$\\sim 12$~m. It is approximately cylindrical in shape with axial and \nforward-backward symmetry about the B0 interaction point. Before going further\nwe describe the coordinate system used at CDF and through this thesis.\n\nB0 is taken as the origin of CDF right-handed coordinate system: $x$-axis\nis horizontal pointing North\\footnote{Outward with respect to the center of \nTevatron.}, $y$-axis is vertical pointing upward and $z$-axis is along beam \nline pointing along proton direction, it identifies {\\em forward} and \n{\\em backward} regions, respectively at $z>0$, East, and $z<0$, West. Sometimes \nit is convenient to work in cylindrical ($r$, $z$, $\\phi$) coordinates where \nthe azimuthal angle $\\phi$ is on the $xy$-plane and is measured from the \n$x$-axis. \nThe $xy$-plane is called \\emph{transverse}, quantities projected on it are \nnoted with a T subscript. Two useful variables are the transverse momentum, $p_{T}$, and energy, $E_{T}$, of a particle:\n\\begin{equation}\n\\vec{p}_{T}\\equiv p\\sin(\\theta)\\mathrm{,}\\qquad E_{T}\\equiv E \\sin(\\theta)\\mathrm{,}\n\\end{equation}\nin collider physics another widely used variable, used in place of $\\theta$, is the \\emph{pseudorapidity}: \n\\begin{equation}\n\\eta\\equiv-\\ln\\big(\\tan (\\theta\/2)\\big).\n\\end{equation}\nIf $(E,\\vec{p})$ is the 4-momentum of a particle, the pseudorapidity is the \nhigh energy approximation ($p\\gg m$) of the \\emph{rapidity}:\n\\begin{equation}\ny=\\frac{1}{2}\\ln\\frac{E+p\\cos(\\theta)}{E-p\\cos(\\theta)}\\stackrel{p \\gg m}{\\rightarrow} \\frac{1}{2}\\ln\\frac{p+p\\cos(\\theta)}{p-p\\cos(\\theta)}= -\\ln\\big( \\tan (\\theta\/2)\\big) \\equiv \\eta.\n\\end{equation}\nA Lorentz boost along the $\\hat{z}$ direction adds a constant $\\ln(\\gamma + \\gamma\\beta)$ to $y$, therefore rapidity differences are invariant.\nThe statistical distribution of final state particles is roughly flat in $y$ because, in hadronic colliders, the interactions between the (anti)proton \nconstituents, which carry only a fraction of the nucleon energy, may have large momentum imbalances along $\\hat{z}$.\n\nFigure~\\ref{CDF_schema} shows an isometric view of the CDF detector and of \nits various sub-detectors. The part inside the $1.4$~T superconducting \nsolenoid contains the integrated tracking system: three silicon sub-detectors \n(the \\emph{Layer00}, the \\emph{Silicon Vertex detector II} and the \n\\emph{Intermediate Silicon Layers}) are the inner core of CDF II. The high \nresolution capability of silicon microstrips is necessary to have good track \nresolution near the interaction point, where particle density is higher. Afterward \nan open cell drift chamber (the \\emph{Central Outer Tracker}) covers until \n$r\\simeq 130$~cm, in the region $|\\eta|<1.0$, the extended lever arm provides \nvery good momentum measurement (\\mbox{$\\Delta p_T\/p^2_T \\simeq 10^{-3}~$GeV$\/c^{-1}$}).\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{Detector\/Figs\/CDF_schema2}\n\\caption[Isometric View of the CDF II Detector]{Isometric view of the CDF II detector, the various sub-detectors are highlighted in different colors and listed.}\\label{CDF_schema}\n\\end{center}\n\\end{figure}\n\nCalorimeter systems are located outside the superconducting solenoid. They are\nbased on \\emph{shower sampling calorimeters} made of sequential layers of \nhigh-Z passive absorbers and active signal generator plastic scintillators. \nThe system is composed by towers with $\\eta - \\phi$ segmentation, each one \ndivided in electromagnetic and hadronic part, they cover the region up to \n$|\\eta|\\simeq 3.6$ ($\\theta\\simeq3^{\\circ}$) and are organized in two main\nsections: the \\emph{Central Calorimeter} covering the region \n$|\\eta|\\lesssim 1.1$ and the \\emph{Plug Calorimeter} \nextending the coverage up to $|\\eta|\\simeq 3.6$. While the central calorimeter\nis unchanged since 1985, the plug calorimeter active part was completely \nrebuilt for Run II, replacing gas chambers with plastic scintillator tiles to\nbetter cope with the higher luminosity.\n\nThe outermost part of CDF detector, outside calorimeters, is occupied by the\nmuon detectors. They are multiple layers of drift chambers arranged in various \nsubsections which cover the region $|\\eta| \\lesssim 1.5$. Only high penetrating\ncharged particles, such as muons, can go across the entire detector.\n\nOther detectors are used for a better particle identification, calibration or \nmonitoring. However a detailed description of the entire CDF detector is far \nfrom the scope of this work. The next paragraphs will focus on tracking and \ncalorimeter systems which play a significant role in the analysis. A complete \ndescription of CDF II detector can be found in~\\cite{TDR}.\n\n\n\n\n\\subsection{Integrated Tracking System}\\label{track_par}\n\nThe trajectory of a charged particle in a uniform \nmagnetic field in vacuum is a helix. A tracking detector identifies some points along particle path so \nthat it is possible to obtain momentum measurements by reconstructing the \nhelix parameters\\footnote{See Section~\\ref{sec:track} for track reconstruction \ndetails.}. A schematic view of CDF tracking volume can be seen in \nFigure~\\ref{CDF_track}: the three main components are the superconducting \nmagnet, the silicon sub-detectors and the central drift chamber.\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.75\\textwidth]{Detector\/Figs\/CDF_track}\n\\caption[CDF II Tracking Volume and Calorimeter]{View of CDF II tracking volume and calorimeter location.}\\label{CDF_track}\n\\end{center}\n\\end{figure}\n\nThe solenoidal magnet, made by NbTi\/Cu superconducting coils, maintains a \nbending magnetic field with a central value of $1.4116$~Tesla, oriented along \nthe positive $\\hat{z}$ direction and nearly uniform in all the tracking volume \n($r\\lesssim 150$~cm and $|z|\\lesssim 250$~cm). \nThe momentum threshold for a particle to radially escape the magnetic field is \n$p_{T}\\gtrsim 0.3$~GeV\/$c$ and the radial thickness of the coil is $0.85$ \nradiation lengths ($X_{0}$).\n\n\n\n\\subsection*{Silicon System}\n\nThe silicon system is the first tracking sub-detector encountered by particles\nexiting from the primary interaction vertex. Semiconductor detectors offer\nexcellent spatial resolution and fast response time.\nTherefore it permits the reconstruction of secondary vertices displaced from the\nprimary, produced in the decay of long lived $b$-hadrons\\footnote{Correct \nidentification of $b$-hadrons is fundamental in many analyses e.g. $b$-hadrons \nare one of the decay products of $top$ quark and also Higgs boson has a high \nbranching ratio to $b$ quarks for $m_H\\lesssim 140$~GeV$\/c^{2}$.}.\n\nCDF employs $\\sim7$~m$^2$ silicon active-surface for a total \nof 722,432 different channels read by about 5500 integrated custom chips. The\ncomplete silicon tracking detector is displayed in Figure~\\ref{silicons}.\nOf the three subsystems composing the core of CDF, the \n\\emph{Layer00 }~\\cite{l00} (L00 ) is the innermost.\nIt consists of a single layer of single-sided silicon sensors directly mounted \non the beam pipe at radii, alternating in $\\phi$, of $1.35$~cm or $1.62$~cm, \ncovering the region $|z|\\lesssim 47$~cm. During the construction of the SVX II \nmicrovertex (see below) CDF realized that the multiple scattering due to the \npresence of read-out electronics and cooling systems installed \ninside tracking volume was going to degrade the impact parameter resolution.\nL\\O\\O~ was designed to recover it thanks to its proximity to the beam. \nFurthermore, being made of state-of-the-art radiation-tolerant sensors, it will\nensure a longer operating lifetime to the entire system.\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.49\\textwidth]{Detector\/Figs\/cdfiisi_rz}\n\\includegraphics[width=0.49\\textwidth]{Detector\/Figs\/cdf_silicon_endview}\n\\caption[Side and Front View of Silicon System]{Side and front view of silicon tracking system at CDF.}\\label{silicons}\n\\end{center}\n\\end{figure}\n\nThe main component of the silicon system is SVX II~\\cite{svx}, the \n\\emph{Silicon VerteX detector}\nis made of three cylindrical barrels for a total length of about $96$~cm along\n$z$, covering the luminosity region until $\\simeq 2.5~\\sigma_{l}$, and with a \npseudo-rapidity\nrange $|\\eta|\\lesssim 2$. Each barrel is divided in twelve identical wedges in\n$\\phi$, arranged in five concentric layers between radii \n$2.4$~cm and $10.7$~cm. Each layer is divided into independent longitudinal \nread-out units, called {\\em ladders}. Each ladder consists of a low-mass support\nfor a double-sided silicon microstrip detector. Three out of five layers \ncombine an $r-\\phi$ measurement on one side with $90^{\\circ}$ stereo measurement\non the other, the remaining two layers combine an $r-\\phi$ measure with a small\nangle $r-z$ stereo measurement (with tilt angle of $1.2^{\\circ}$). \nThe highly parallel fiber based data acquisition \nsystem reads out the entire sub-detector in approximately $10~\\mu$s.\n\nThe \\emph{Intermediate Silicon Layers} detector~\\cite{isl} (ISL) is the outermost of the three silicon sub-detectors, radially\nlocated between SVX II and the drift chamber covering the region $|\\eta|\\lesssim 2$. It is divided in three barrels segmented into $\\phi$ \nwedges. The central barrel ($|\\eta|\\lesssim 1$) is made of one layer of silicon sensors at radius of $22$~cm, instead the two outer barrels \n($1\\lesssim|\\eta|\\lesssim 2$) are made of two layers at radii of $20$~cm and \n$28$~cm. Its purpose is to strengthen the CDF tracking in the central region and to add precision hits in a region not fully covered by the \ndrift chamber. Track reconstruction can be extended to the whole region $|\\eta|<2$ using the silicon detector.\n\nThe complete silicon sub-detector (L00, SVX II and ISL) has an asymptotic \nresolution of $40~\\mu$m in impact parameter and of $70~\\mu$m along $z$ \ndirection. The total amount of material varies roughly as:\n\\begin{equation}\n\\frac{0.1X_{0}}{\\sin(\\theta)}\n\\end{equation}\nin the central region and doubles in the forward region because of the presence of read-out electronics, cooling system and support frames~\\cite{sil_performance}.\n\n\n\n\\subsection{Central Outer Tracker}\n\nThe \\emph{Central Outer Tracker}~\\cite{cot} (COT) is an open-cell drift \nchamber used for particles tracking at large radii. It has an \nhollow-cylindrical geometry and covers $43.3 10$ are removed. The procedure is iterated until all accepted tracks have $\\chi^2< 10$.\nA quality index (see Table~\\ref{tab:pVtx}) is assigned to the primary vertex depending on track multiplicity and type: a quality $\\ge 12$ is required for primary interaction vertex reconstruction. \n\nThe PV position is defined by $\\left(x_{PV},y_{PV},z_{PV}\\right)$. Typical $x_{PV}$ and $y_{PV}$ are of the order of tens of microns while a cut of $|z_{PV}|\\leq 60$~cm (luminosity region fiducial selection) is applied to constrain the collisions in the geometrical region where the detector provides optimal coverage.\n\n\n \\begin{table}[h]\n\\begin{center}\n \\begin{tabular}{cc}\\toprule\n Criterion & Quality Value\\\\\\midrule\n Number Si -tracks$\\ge$3 & 1\\\\\n Number Si -tracks$\\ge$6 & 3\\\\\n Number COT-tracks$\\ge$1 & 4\\\\\n Number COT-tracks$\\ge$2 & 12\\\\\n Number COT-tracks$\\ge$4 & 28\\\\\n Number COT-tracks$\\ge$6 & 60\\\\\n \\bottomrule\n \\end{tabular}\n \\caption[Primary Vertex Quality Criteria]{Primary Vertex quality criteria: a quality $\\ge 12$ is required for primary interaction vertex reconstruction.}\\label{tab:pVtx}\n\\end{center}\n \\end{table}\n\n\n\n\n\\section{Lepton Identification Algorithms}\\label{sec:lep_id}\n\nIn this section the specific lepton identification algorithms\\footnote{A lepton is, by the experimental point of view, an electron or a $\\mu$ with no distinction between a particle and its anti-particle. Also $\\tau$ leptonic decays can enter in the lepton sample but the algorithms are not optimized for them.} are discussed. We distinguish two {\\em tight electron} identification algorithms (CEM, PHX), two {\\em tight muon} identification algorithms (CMUP, CMX), six {\\em loose muon} identification algorithms (BMU, CMU, CMP, SCMIO, CMIO, CMX-NotTrig) and one isolated track identification algorithm (ISOTRK). \nFigure~\\ref{fig:split_alg} shows, for a $WZ$ MC, the detector $\\eta-\\phi$ coverage separately for each lepton identification category. Figure~\\ref{fig:all_alg} shows them for all the categories together. The eleven lepton identification algorithms are described in the following sections.\n\n\\begin{figure}[h!]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{PhysObj\/Figs\/lepSet1}\n\\includegraphics[width=0.495\\textwidth]{PhysObj\/Figs\/lepSet2}\\\\\n\\includegraphics[width=0.495\\textwidth]{PhysObj\/Figs\/lepSet3}\n\\includegraphics[width=0.495\\textwidth]{PhysObj\/Figs\/lepSet4}\n\\caption[Coverage, in the $\\eta - \\phi$ Plane, for the Lepton Identification Algorithms]\n {Coverage of the detector $\\eta - \\phi$ plane for the different lepton identification algorithms evaluated on a $WZ$ Monte Carlo: CEM and PHX tight electrons (top left); CMUP and CMX tight muons (top right); BMU, CMU, CMP, SCMIO, CMIO, CMX-NotTrig loose muons (bottom left); ISOTRK isolated tracks (bottom right).}\\label{fig:split_alg}\n\\end{center}\n\\end{figure}\n\n\\begin{figure}[h!]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{PhysObj\/Figs\/lepSet0}\n\\caption[Combined coverage, in the $\\eta - \\phi$ Plane, for the Lepton Identification Algorithms Combined]\n{Coverage of the detector $\\eta - \\phi$ plane for all the combined lepton identification algorithms evaluated on a $WZ$ Monte Carlo.}\\label{fig:all_alg}\n\\end{center}\n\\end{figure}\n\n\n\n\\subsection{Electron Identification}\\label{sec:eleSel}\n\nA candidate electron is ideally defined as an energy deposit ({\\em cluster}) in the EM section of the calorimeters and a charged track matched to it. The track requirements removes most of the ambiguity due to photon EM showers. \n\nThe two electron identification algorithms used in this analysis are divided according to the CDF calorimeter segmentation in CEM (Central ElectroMagnetic object) and PHX (Plug electromagnetic object with pHoeniX tracking), for $|\\eta|<1.1$ and $1.2<|\\eta|<3.6$ respectively.\n\nThe CDF EM clustering algorithm~\\cite{ele_cal} works in a simple but efficient way. The physical space corresponding to the calorimeter towers is mapped in the $\\eta-\\phi$ plane and the algorithm creates two lists of the calorimeter towers ordered by decreasing energy measurement: the {\\em usable list} (working towers with energy $>100$~MeV) and the {\\em seed list} (towers with energy $>2$~GeV). Then, it takes the first seed tower and creates an $\\eta-\\phi$ cluster by adding the neighboring towers to form a $2\\times 2$ or $3\\times 3$ $\\eta-\\phi$ area. An EM cluster is found if:\n\\begin{equation}\n E^{Had}\/E^{EM}<0.125\\textrm{,}\n\\end{equation}\nwhere $E^{Had}$ is the energy deposited in the hadronic calorimeter section and $E^{EM}$ is the corresponding quantity for the EM section. As final step, the $\\eta-\\phi$ centroid of the cluster is calculated and the used towers are removed from the list. The algorithm selects the next seed tower and iterates the process until there are no more seed towers available. \n\nThe $3\\times 3$ clustering is used for the CEM algorithm while the $2\\times 2$ clusters are used in the Plug region. A cluster is not allowed to cross the boundary between different sub-detectors. \n\nSeveral corrections are applied to reconstruct the final energy: lateral leakage, location inside the physical tower, on-line calibration and response curve drawn by the test beam data. Also the energy measured in the shower-max (CES, PES) and preshower (CPR, PPR) detectors is added to the final reconstructed energy. The shower-max profile is also compared to the calibration profiles of electrons or photons and, last but not least, it is used to measure the position of the EM shower centroid. \n\nBeyond the EM energy measurement, the calorimeter information is further exploited for a better particle identification. The following variables are used:\n\\begin{itemize}\n\\item $E^{Had}\/E^{EM}$ ratio: studies performed with candidate $Z^0\\to e^+e^-$ events~\\cite{w_z_prl} show that electrons detected in the central or in the plug region have a little deposit in the hadronic part of the calorimeter (Figure~\\ref{ele_deposit}). \n\\item {\\em Lateral shower sharing} variable, $L_{shr}$, compares the sharing of energy deposition between the towers in the CEM to the expected in true electromagnetic showers taken with test beam data:\n\\begin{equation}\n L_{shr}=0.14\\sum_i \\frac{E_i^{adj}- E_i^{expect}}{\\sqrt{(0.14\\sqrt{E_i^{adj}})^2+(\\Delta E_i^{expect})^2}}\\mathrm{,}\n\\end{equation}\nwhere the sum is over the towers adjacent (\\emph{adj}) to the seed tower of the cluster, $0.14\\sqrt{E_i^{adj}}$ is the error on the energy measure and $\\Delta E_i^{expect}$ is the error on the energy estimate.\n\\item The $\\chi^2$ of the fit between the energy deposit and the one obtained from test beam data ($\\chi^2_{strip}$ for CEM and $\\chi^2_{towers}$ for PHX).\n\\end{itemize}\n\\begin{figure}[!h]\n\\begin{center}\n\\includegraphics[width=0.92\\textwidth]{PhysObj\/Figs\/ele_deposit_cen}\\\\\n\\includegraphics[width=0.445\\textwidth]{PhysObj\/Figs\/ele_hadem_plug}\n\\includegraphics[width=0.445\\textwidth]{PhysObj\/Figs\/ele_iso_plug}\n\\caption[Electron Identification Criteria for $Z^0\\to e^+e^-$ Candidate Events]{$E^{Had}\/E^{EM}$ (left) and isolation (right) distribution of central (top) and plug (bottom) calorimeter electron selection from unbiased, second legs of $Z^0\\to e^+e^-$ candidate events in data~\\cite{w_z_prl}.}\\label{ele_deposit}\n\\end{center}\n\\end{figure}\n\nFinally, a reconstructed track matched to the EM cluster is used to suppress photons in the central region. Due to poor tracking in the forward region, PHX electrons are defined by a different strategy that relies on the {\\em Phoenix} matching scheme~\\cite{phx_alg}. \n\nTracks candidate electrons in the central region must satisfy the following requirements:\n\\begin{itemize}\n\\item track quality: $\\ge 3$ ($\\ge 2$) COT Axial (Stereo) segments with at least five hits each associated with the track.\n\\item $E^{cluster}\/p^{trk}<2$: bremsstrahlung photons are emitted colinearly to the electron and energy is radiated in the same EM cluster, therefore the measured track momentum is lower than the original $p_T$ of the electron.\n\\item Track-shower matching: \n \\begin{equation}\n 3.01.1$. However to provide some fake-electron rejection a road is built starting from the $x, y$ position in the PES detector, the PV position and the curvature provided by the $E_T$ cluster. The Phoenix matching succeeds if at at last three hits in the silicon detectors are found, thus allowing the selection of a PHX electron candidate. This is not considered a {\\em real} track matching because only four points are used in the fit and it is impossible to reconstruct a five-parameter helix, moreover some events have a high density of silicon hits, increasing the ambiguity of the matching. \n\nThe other important identification requirement is the {\\em Isolation} or $IsoRel$, a variable describing how much calorimeter activity surrounds the lepton. It is defined as: \n\\begin{equation}\n IsoRel\\equiv E_T^{iso}\/E_T^{cluster}<0.1,\\quad E_T^{iso}=E_T^{0.4}-E_T^{cluster},\n\\end{equation}\n$E_T^{0.4}$ is the energy collected by the calorimeters within a radius $\\Delta R=0.4$ from the centroid of the EM cluster. Isolation is used in analyses involving a $W^{\\pm}$ or $Z^0$ boson because the kinematic region allowed to leptons coming from the bosons decay is usually far from jets or other particles (see Figure~\\ref{ele_deposit}).\n\nTable~\\ref{tab:ele_alg} summarizes all the CEM and PHX identification criteria while the top left part of Figure~\\ref{fig:split_alg} shows the detector coverage of the two algorithms in the $\\eta - \\phi$ plane for a $WZ$ Monte Carlo.\n\\begin{table}\\begin{center}\n\\begin{tabular}{cc}\n\\toprule\nElectrons & Identification Cuts\\\\\n\\midrule\n& EM fiduciality \\\\\n& $E_T^{EM}>20$~GeV \\\\\n& $E^{Had}\/E^{EM}<0.055+0.0045 E^{EM}$ \\\\\n& $L_{shr}<0.2$ \\\\\n& $\\chi^2_{\\mathrm{{CES strip}}}<10$ \\\\\nCEM &$\\ge 3$ axial $\\ge 2$ stereo COT Layers \\\\\n& $\\Delta z(CES,trk)<3$~cm \\\\\n& $-3.0100$~GeV \\\\\n&$|z_0^{trk}-z_0^{PV}|<5$~cm \\\\\n& $|z_0^{trk}|<60$~cm \\\\\n&$IsoRel<0.1$ \\\\\n\\midrule\n& $1.2<|\\eta^{PES}|<2.8$ \\\\\n& $E_T^{EM}>20$~GeV \\\\\n& $E^{Had}\/E^{EM}<0.05$ \\\\\n& PEM$3\\times 3$Fit $\\ne 0$ \\\\\n& $\\chi^2_{PEM3\\times 3} <10$ \\\\\nPHX & $PES5\\mathrm{by}9U>0.65$ \\\\\n& $PES5\\mathrm{by}9V>0.65$ \\\\\n& $\\Delta R(PES,\\mathrm{cluster})<3$~cm \\\\\n& N Si Hits $\\ge 3$ \\\\ \n& $\\mathrm{Phoenix~Track}$\\\\\n&$|z_0^{PHX}-z_0^{PV}|<5$~cm \\\\\n& $|z_0^{PHX}|<60$~cm \\\\\n&$IsoRel<0.1$ \\\\\n\\bottomrule\n\\end{tabular}\n\\caption[CEM and PHX Electron Selection Requirements]{Summary of CEM and PHX \\emph{tight electron} selection requirements.}\\label{tab:ele_alg}\n\\end{center}\\end{table}\n\n\n\n\\subsection{Muon Identification}\n\nA muon behaves like a minimum ionizing particle due to its rest mass, which is about 200 times larger than the electron one. Therefore muons deposit very little energy in the calorimeter systems and can leave a signal in the outer layer of the CDF detector which is instrumented with arrays of gas detector and scintillators (muon chambers). Ionization deposits from a muon candidate in a given muon detector constitute a {\\em stub} and the candidate naming depends by the coverage of the muon detector that records them (see Section~\\ref{sec:mu_cham}). \n\nThe basic selection for a muon candidate~\\cite{muon_id} is a high quality COT track pointing to a Minimum Ionizing Particle (MIP) energy deposit in the EM and HAD calorimeters and matched to a stub in the muon chambers. The precise requirements are the following:\n\\begin{itemize}\n\\item track quality: the reconstructed COT track must have a minimal amount of axial ($\\ge 3$) and stereo ($\\ge 2$) COT super-layers;\n\\item the $\\chi^2_{trk}$ returned by the track fitting algorithm should be less than $2.3$;\n\\item track\/stub matching in the $\\phi$-plane is required by an appropriate $\\Delta x(trk,\\mathrm{stub}) = | x^{trk} - x^{\\mathrm{stub}}|$ with cut values different for each specific muon sub-detector\n\\end{itemize}\n\nIf we are interested to a muon coming from $Z$ or $W$ decay, we expect the muon to be isolated from other detector activity (as in the electron case). The isolation is defined exploiting the muon candidate momentum and MIP energy deposit in the calorimeter:\n\\begin{equation}\n IsoRel\\equiv E_T^{iso}\/p_T<0.1;\\quad E_T^{iso}=E_T^{0.4}-E_T^{MIP},\n\\end{equation}\nwhere the $p_T$ is the COT track momentum and $E_T^{MIP}$ is the transverse energy deposited in the towers crossed by the track.\n\nThe two {\\em tight} muon reconstruction algorithms are the CMUP and CMX. The first covers the region $|\\eta|<0.6$ where a track is required to match stubs in both CMU and CMP muon detectors. The second covers the region $0.65<|\\eta|<1.0$: it requires a stub in the CMX muon detector and a minimum curvature of the COT track, $\\rho_{COT}$, of $140$~cm\\footnote{This last requirement ensures appropriate efficiency of the CMX trigger.}. Table~\\ref{tab:central_mu_alg} summarizes all the CMUP and CMX identification criteria while the top right section of Figure~\\ref{fig:split_alg} shows the detector coverage of the two algorithms in the $\\eta - \\phi$ plane for a $WZ$ Monte Carlo.\n\\begin{table}\\begin{center}\n\\begin{tabular}{cc}\n\\toprule\nMuons & Identification Cuts\\\\\n\\midrule\n&$p_T^{trk}>20$~GeV$\/c$ \\\\\n&$E_T^{EM}<2+\\mathrm{max}\\big(0, ((p^{trk}-100)0.0115)\\big)$~GeV \\\\\n&$E_T^{Had}<6+\\mathrm{max}\\big(0, ((p^{trk}-100)0.028)\\big)$~GeV \\\\\n& $\\ge 3$ axial $\\ge2$ stereo COT Layers \\\\\nCMUP & $|d_0^{trk}|<0.2$~cm (track w\/o silicon) \\\\\nCMX & $|d_0^{trk}|<0.02$~cm (track with silicon) \\\\\n& $|z_0^{trk}-z_0^{PV}|<5$~cm \\\\\n& $|z_0^{trk}|<60$~cm \\\\\n& $\\chi^2_{trk}<2.3$\\\\\n&$IsoRel<0.1$ \\\\ \\midrule\n& CMU Fiduciality \\\\ \nCMUP & CMP Fiduciality \\\\\n& $\\Delta X_{CMU}(trk, \\mathrm{stub})<7$~cm \\\\\n& $\\Delta X_{CMP}(trk, \\mathrm{stub})<5$~cm \\\\\\hline\n& CMX Fiduciality \\\\ \nCMX & $\\rho_{COT}>140$~cm \\\\\n& $\\Delta X_{CMX}(trk, \\mathrm{stub})<6$~cm \\\\\n\\bottomrule\n\\end{tabular}\n\\caption[CMUP and CMX Muon Selection Requirements]{Summary of CMUP and CMX \\emph{tight muons} selection requirements.}\\label{tab:central_mu_alg}\n\\end{center}\\end{table}\n\n\n\n\n\\subsection{Loose Muons Identification}\\label{sec:LooseMuId}\n\nMuons can be faked by cosmic rays or hadrons showering deep inside the calorimeters or not showering at all, however the muon candidates usually have a very clean signature because of the many detector layers used to identify them. On the other hand the signal acceptance for the tight muon categories is geometrically limited, therefore a set of lower quality muon identification criteria, {\\em loose muons}, was developed. Loose muons were used in other analysis~\\cite{single_top} to increase the signal yield. These are the main requirements of the loose muons algorithms:\n\\begin{itemize}\n\\item BMU: forward isolated muons ($1.0<|\\eta|<1.5$) with hits in {\\em Barrel} muon chambers.\n\\item CMU: central isolated muons with hits {\\em only} in the CMU chambers and not in CMP.\n\\item CMP: central isolated muons with hits {\\em only} in the CMP chambers and not in CMU.\n\\item SCMIO: a good quality track matched to an isolated MIP deposit and a {\\em non-fiducial} stub in the muon detector (Stubbed Central Minimum Ionizing particle).\n\\item CMIO: a good quality track {\\em only} matched to an isolated MIP deposit and failing any other muon identification criteria (Central Minimum Ionizing particle).\n\\item CMX-NotTrigger: also named CMXNT, are isolated muon detected in the CMX chamber but which are {\\em not} triggered by the CMX specific triggers because of the geometrical limits of the COT ($\\rho_{COT}<140$~cm).\n\\end{itemize}\nTable~\\ref{tab:loose_alg} summarizes all the Loose muons identification criteria while the bottom left section of Figure~\\ref{fig:split_alg} shows the detector coverage of the six algorithms in the $\\eta - \\phi$ plane for a $WZ$ Monte Carlo.\n\n\n\\begin{table}\\begin{center}\n\\begin{tabular}{cc}\n\\toprule\nMuons & Identification Cuts \\\\\n\\midrule\n&$p_T^{trk}>20$~GeV$\/c$ \\\\\n&$E_T^{EM}<2+\\mathrm{max}\\big(0, ((p^{trk}-100)0.0115)\\big)$~GeV \\\\\nBMU &$E_T^{Had}<6+\\mathrm{max}\\big(0, ((p^{trk}-100)0.028)\\big)$~GeV \\\\\nCMU & $|d_0^{trk}|<0.2$~cm (track w\/o silicon) \\\\\nCMP & $|d_0^{trk}|<0.02$~cm (track with silicon) \\\\\nSCMIO& $|z_0^{trk}-z_0^{PV}|<5$~cm \\\\\nCMIO& $|z_0^{trk}|<60$~cm \\\\\nCMXNT& $\\chi^2_{trk}<2.3$\\\\\n&$IsoRel<0.1$ \\\\ \\midrule\n\n& BMU Fiduciality \\\\ \nBMU & $\\Delta X_{BMU}(trk, \\mathrm{stub})<9$~cm \\\\\n& COT Hits Frac. $>0.6$ \\\\\\midrule\n\n& CMU Fiduciality \\\\ \nCMU & $\\Delta X_{CMU}(trk, \\mathrm{stub})<7$~cm \\\\\n& $\\ge 3$ axial $\\ge 2$ stereo COT Layers \\\\\\midrule\n\n& CMP Fiduciality \\\\ \nCMP & $\\Delta X_{CMP}(trk, \\mathrm{stub})<7$~cm \\\\\n& $\\ge 3$ axial $\\ge 2$ stereo COT Layers \\\\\\midrule\n\n& Stub Not Fiducial \\\\ \nSCMIO & $\\ge 3$ axial $\\ge 2$ stereo COT Layers \\\\\n&$E_T^{EM} + E_T^{Had} > 0.1$ \\\\\\midrule\n\n& No Stub \\\\ \nCMIO & $\\ge 3$ axial $\\ge 2$ stereo COT Layers \\\\\n&$E_T^{EM} + E_T^{Had} > 0.1$ \\\\\\midrule\n\n\n& CMX Fiduciality \\\\ \nCMXNT & $\\Delta X_{CMX}(trk, \\mathrm{stub})<6$~cm \\\\\n& $\\ge 3$ axial $\\ge 2$ stereo COT Layers \\\\\n& $\\rho_{COT}<140$~cm, Not trigger CMX \\\\\n\\bottomrule\n\\end{tabular}\n\\caption[Loose Muons Selection Requirements]{Summary of the BMU, CMU, CMP, SCMIO, CMIO, CMX-Not-Trigger (CMXNT) {\\em Loose muons} selection requirements.}\\label{tab:loose_alg}\n\\end{center}\\end{table}\n\n\n\n\n\n\n\\subsection{Isolated Track Identification}\\label{sec:IsoTrk}\n\nThe isolated tracks are last lepton category used in this analysis. They are defined to be high-$p_T$ good quality tracks isolated from energy deposits in the tracking systems. The track isolation is defined as:\n\\begin{equation}\\label{TrackIsolation}\n\\rm{TrkIso}=\\frac{p_T({\\rm{trk\\ candidate}})}{p_T({\\rm{trk\\ candidate}})+\\sum p_T({\\rm{other\\ trk)}}}\\,\\rm{,}\n\\end{equation}\nwhere $p_T({\\rm{trk\\ candidate}}$ is the transverse momentum of the specific track we analyze (candidate) and $\\sum p_T({\\rm{other\\ trk)}}$ is the sum of the transverse momenta of all good quality tracks within a cone radius of 0.4 of the candidate track. The isolation requirement is necessary in order to ensure that the track corresponds to a charged lepton produced in a decay of a $W$ boson and it is not part of an hadronic jet. A track is fully isolated if $TrkIso = 1.0$, thus a cut value of $0.9$ is used in the analysis.\nA selection requirement of $\\Delta R>0.4$ between the track and any tight jet\\footnote{Central, high $E_T$ jets are classified as {\\em tight}: $|\\eta|<2.0$, $E_T^{cor}>20$ (see Section~\\ref{sec:jetObj}).} is also applied to remove jets with low track multiplicity.\n\nThe fact that the isolated track is not required to match a calorimeter cluster or a muon stub allows to recover real charged leptons that arrive in non-instrumented regions of the calorimeter or muon detectors, as seen in Figure~\\ref{fig:split_alg}. \n\nTable~\\ref{iso_tab} summarizes the criteria used to select good quality isolated tracks while the bottom right section of Figure~\\ref{fig:split_alg} shows the detector coverage of the algorithms in the $\\eta - \\phi$ plane for a $WZ$ Monte Carlo.\n\n\\begin{table}\\begin{center}\n\\begin{tabular}{cc}\n\\toprule\nTrack Leptons & Identification Cuts \\\\\n\\midrule\n&$p_T^{trk}>20$~GeV$\/c$ \\\\\n& $\\ge 24$ axial $\\ge 20$ stereo COT Hits \\\\\n& $|d_0^{trk}|<0.2$~cm (track w\/o silicon) \\\\\n& $|d_0^{trk}|<0.02$~cm (track with silicon)\\\\\nIsolated & $|z_0^{trk}-z_0^{PV}|<5$~cm \\\\\nTracks & $|z_0^{trk}|<60$~cm \\\\\n& $\\chi^2_{trk}$ probability $<10^{-8}$\\\\\n& N Si Hits $\\ge 3$ \\\\ \n&$TrkIso>0.9$ \\\\\n& $\\Delta R(trk, \\mathrm{tight~jets})>0.4$\\\\\n& $\\Delta \\phi(trk, \\mathrm{jet 1})>0.4$\\\\\n\\bottomrule\n\\end{tabular}\n\\caption[Isolated Tracks Selection Requirements]{Summary of ISOTRK \\emph{isolated track} leptons selection requirements.}\\label{iso_tab}\n\\end{center}\\end{table}\n\n\n\n\\section{Jet Identification}\\label{sec:jetObj}\n\nThe QCD theory tells us that the partons composing the (anti)proton can be treated perturbatively as free particles if they are stuck by an external \nprobe\\footnote{I.e. a lepton, a photon or a parton from a different hadron.} with sufficient high energy (so called {\\em hard scattering}). However partons resulting from the \ninteraction can not exist as free particles because at longer distances (i.e. lower energies) the \\emph{strong potential} can not be treated perturbatively and \npartons must form colorless hadrons. This process, called {\\em hadronization} or {\\em showering}, produces a {\\em jet}, i.e. a narrow spray of stable particles that retains the information of the initial parton (for a pictorial representation see Figure~\\ref{jet1}). \n\\begin{figure}[t]\n\\begin{center}\n\\includegraphics[width=0.6\\textwidth]{PhysObj\/Figs\/jet1}\n\\caption[Schematic Representation of a Hadronic Jet]{A parton originating from a hard scattering hadronizes and generates a\na narrow spray of particles identified as a {\\em jet}.}\\label{jet1}\n\\end{center}\n\\end{figure}\n\nFrom an experimenter's point of view a jet is defined as a large energy deposit in a localized area of the detector (see Figure~\\ref{lego_jet}). The challenge of a physics analysis\nis to recover from detector information the initial energy, momentum and, possibly, the kind of the parton produced in the original interaction. \n\\begin{figure}[t]\n\\begin{center}\n\\includegraphics[width=0.85\\textwidth]{PhysObj\/Figs\/lego_jet}\n\\caption[Event Display with Jet Deposits]{Calorimetric deposit in the $\\eta-\\phi$ plane as represented in the CDF event display. EM deposits are red while HAD deposits are blue.}\\label{lego_jet}\n\\end{center}\n\\end{figure}\nA {\\em jet identification algorithm} is a tool to reconstruct such information and it should satisfy at best the following requirements~\\cite{jet_req}:\n\\begin{itemize}\n\\item\\emph{Infrared safety}: the presence of soft radiation between two jets may cause a merging of the two jets. This should not occur to avoid an incorrect parton attribution.\n\\item\\emph{Collinear safety}: the jet reconstruction should be independent from any collinear radiation in the event, i.e. different energy distribution of particles inside calorimetric towers.\n\\item\\emph{Invariance under boost}: the same jets should be found independently from boosts in longitudinal direction.\n\\item\\emph{Boundary stability}: kinematic variables should be independent from the details of the final state.\n\\item\\emph{Order independence}: the same reconstructed quantities should appear looking at parton, particle and detector levels.\n\\item\\emph{Straightforward implementation}: algorithm should be easy to implement in perturbative calculations.\n\\end{itemize}\nBeyond this theoretical aspects a jet algorithm should be experimentally efficient with a high reconstruction efficiency, good resolution and robust at high instantaneous luminosity.\n\n\n\\subsection{CDF Cone Algorithm}\\label{sec:jetCone}\n\nCDF uses several algorithms, none of them completely satisfying all the above requirements. The most common one, that is also the one used in this ana\\-lysis,\nis \\verb'JETCLU'~\\cite{jetclu}, an iterative fixed cone jet reconstruction algorithm based only on calorimetric information.\n\nThe algorithm starts by creating a list of the seed towers from all the calorimeter towers with transverse energy above the threshold of $1$~GeV. \nStarting with the highest-$E_T$ seed tower, a precluster is formed by combining\ntogether all adjacent seed towers within a cone of given radius $R$ in the $\\eta-\\phi$ space\\footnote{In this analysis we use $R=0.4$.}. This procedure is repeated, starting with the next unused seed tower, until the list is exhausted. The $E_T$-weighted centroid is then formed from the towers in the precluster and a new cone of radius $R$ is formed around this centroid. All towers with energy above the lower threshold of $100$~MeV within this new cone are added to the cluster. Then, a new centroid is calculated from the set of towers within the cluster and a new cone drawn. This process is iterated until the centroid of the energy deposition within the cone is aligned with the geometric axis of the cone (stable solution).\n\nThe initial clusters found can overlap so the next step is to merge or separate overlapping clusters, since each tower may belong only to one jet, each particle should not be assigned to more than one jet. Two clusters are merged if the total energy of the overlapping towers is greater than $75\\%$ of the energy of the smallest cluster. If the shared energy is below this cut, the shared towers are assigned to the cluster that is closer in $\\eta-\\phi$ space. The process is iterated again until the list of clusters remains fixed. \n\nThe final step of the jet identification happens on a second stage, after electron candidate identification and it is called {\\em reclustering}. If an EM calorimeter cluster is found to be compatible with a tight electron identification (CEM, PHX both isolated or not), the EM calorimeter towers are removed and the jet clustering algorithm is iterated.\n\nMassless four-vector momenta are assigned to the towers in the clusters for EM and HAD components with a magnitude equal to the energy deposited in the tower and the direction defined by a unit vector pointing from the event vertex to the center of the calorimeter tower at the depth that corresponds to the shower maximum. A cluster four-vector is then defined summing over the towers in the cluster:\n\\begin{eqnarray}\n&& E=\\sum_{i=1}^N(E_i^{EM}+E_i^{HAD})\\mathrm{,}\\\\\n&& p_x=\\sum_{i=1}^N(E_i^{EM}\\sin\\theta_i^{EM}\\cos\\phi_i^{EM}+E_i^{HAD}\\sin\\theta_i^{HAD}\\cos\\phi_i^{HAD})\\mathrm{,}\\\\\n&& p_y=\\sum_{i=1}^N(E_i^{EM}\\sin\\theta_i^{EM}\\sin\\phi_i^{EM}+E_i^{HAD}\\sin\\theta_i^{HAD}\\sin\\phi_i^{HAD})\\mathrm{,}\\\\\n&& p_z=\\sum_{i=1}^N(E_i^{EM}\\cos\\theta_i^{EM}+E_i^{HAD}\\cos\\theta_i^{HAD})\\mathrm{.}\n\\end{eqnarray}\nwhere the index $i$ runs over the towers in the cluster.\n\nIn order to study jet characteristics, other variables (number of tracks, energy deposited in the HAD and EM calorimeters, etc.) are reconstructed and associated to the final jet analysis-object.\n\n\n\n\\subsection{Jet Energy Corrections}\n\nThe ultimate goal of the jet reconstruction algorithm is the determination of the exact energy of the outgoing partons coming from the hard interaction, i.e. the Jet Energy Scale (JES). Clearly many factors produce a mismatch between the raw energy measured by the algorithm and the one of the parton before the hadronization. \n\nCDF developed a set of jet energy corrections depending of $\\eta$, $E_T^{raw}$ and $R$ of the jet reconstructed by \\verb'JETCLU' algorithm. The corrections are divided into five levels\\footnote{The actual naming skips $L2$, because it is absorbed in $L1$, and $L3$, as it was introduced as a temporary MC calibration in Run II.} ($L$-levels) that can be applied in a standard way to different analyses~\\cite{jet_corr1, jet_corr2}: $\\eta$-dependent response($L1$), effect of multiple interactions ($L4$), absolute energy scale ($L5$), underlying event ($L6$) and out-of-cone ($L7$) corrections. The correction $L1$ and $L5$ are multiplicative factors ($f_{L1}$ and $f_{L5}$) on the raw $E_T$ of the jet, the others are additive constants ($A_{L4}$, $A_{L6}$ and $A_{L7}$). The equation for the complete correction is:\n\\begin{equation}\\label{eq:fullCor}\nE_T^{FullCor}(\\eta, E_T^{raw}, R)=(E_T^{raw} f_{L1}-A_{L4})f_{L5}-A_{L6}+A_{L7}.\n\\end{equation}\nA description of each term is given in the following paragraphs while Figure~\\ref{totsys} shows the separate systematic uncertainties ($\\sigma_{JES}$) associated to each term of Equation~\\ref{eq:fullCor}.\n\nRecent studies~\\cite{gluonJES_10829} have shown that the simulated detector response to jets originating from high-$p_T$ gluons is improved by lowering the reconstructed $E_T$ by two times the uncertainty used in Equation~\\ref{eq:fullCor} ($-2\\sigma_{JES}$). More accurate studies are ongoing within the CDF collaboration to fully understand the effect of a {\\em parton dependent} jet energy correction. The $-2\\sigma_{JES}$ gluon-jets prescription has been applied also to this analysis, however the result is negligible due to the small fraction of gluon-jets present in the $HF$-enriched signal region.\n\n\\begin{figure}[t]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{PhysObj\/Figs\/totsys2}\n\\caption[Jet Energy Scale Systematic Uncertainty]{Jet $p_T$ dependence of the systematic uncertainties ($\\sigma_{JES}$) relative to the different terms of Equation~\\ref{eq:fullCor} for the \\texttt{JETCLU} algorithm with $R=0.4$ cone size.}\\label{totsys}\n\\end{center}\n\\end{figure}\n\nIn this analysis, like in many others~\\cite{single_top, wh94_note10796}, we choose to use jet corrections only up to L5. Therefore we define the Level-5 only correction:\n\\begin{equation}\nE_T^{cor}(\\eta, E_T^{raw}, R)=(E_T^{raw} f_{L1}-A_{L4})f_{L5}.\n\\end{equation}\n\nL6 and L7 corrections are fundamental for the measurement of the unknown mass of a particle in the hadronic final state (for example top quark mass measurements) but they are much less relevant in a cross section measurement where only the relative data\/MC energy scale matters and not its absolute value. Furthermore, Figure~\\ref{totsys} shows that a large uncertainty is associated to the out-of-cone correction for $R=0.4$ jets, which would decrease the di-jet invariant mass resolution, lowering the sensitivity of the present measurement.\n\nDepending on, L5 corrected, $E_T^{cor}$ and the jet centroid position in the detector, jets are classified as {\\em tight} or {\\em loose}:\n\\begin{itemize}\n\\item{\\bf Tight:} $|\\eta^{det}|<2.0$, $E_T^{cor}>20$~GeV.\n\\item{\\bf Loose:} $|\\eta^{det}|<2.4$, $E_T^{cor}>12$~GeV and the jet is not tight.\n\\end{itemize}\nEvent selection (see Section~\\ref{sec:jetSel}) is based on the number of tight jets while both tight and loose jets are used in the \\mbox{${\\not}{E_T}$} correction. Jet clusters of lower energy are not analyzed, it is assumed that they produce a negligible noise on the top of the unclustered energy of the event.\n\n\n\n\n\n\\paragraph{Level-1: $\\eta$ Dependent Corrections}\n\n$L1$ correction is applied to raw jet energy measured in the calorimeters to \nmake the detector response uniform in $\\eta$, it takes into account aging of\nthe sub-detectors\\footnote{This was the $L2$ correction during Run I} and other \nhardware non-uniformities (for example the presence of cracks).\nThis correction is obtained using a large di-jet sample:\nevents with one jet (\\emph{trigger jet}) in the central region of the \ncalorimeter ($0.2<|\\eta|<0.6$), where the detector response is well known and \nflat in $\\eta$, and a second jet (\\emph{probe jet}), allowed to range anywhere\nin the calorimeter ($|\\eta|<3.6$). In a perfect detector the jets should be\nbalanced in $p_T$, a balancing fraction is formed:\n\\begin{equation}\nf_b\\equiv \\frac{\\Delta p_T}{p_T^{ave}}=\\frac{p_T^{probe}-p_T^{trigger}}{(p_T^{probe}+p_T^{trigger})\/2}\\mathrm{,}\n\\end{equation}\nthe average of $f_b$ in the analyzed $\\eta$ bin is used to define the $\\beta$\nfactor\\footnote{The definition of Equation~\\ref{beta_eq} has a average value equal\nto $p_T^{probe}\/p_T^{trig}$ but is less sensitive to presence of non-Gaussian \ntails in the usual $p_T^{probe}\/p_T^{trig}$ ratio.} (Figure~\\ref{eta_correction} \nshows the $\\beta$ distribution for different cone radii):\n\\begin{equation}\\label{beta_eq}\n\\beta\\equiv\\frac{2+}{2-}\\mathrm{.}\n\\end{equation}\nThe final $L1$ correction is defined as $f_{L1}(\\eta, E_T^{raw}, R)=1\/\\beta$ \nand reproduces an approximately flat response in $\\eta$ with an error varying \nfrom $0.5\\%$ to $7.5\\%$.\n\\begin{figure}[t]\n\\begin{center}\n\\includegraphics[width=0.85\\textwidth]{PhysObj\/Figs\/eta_correction}\n\\caption[Jet Energy Correction as a Function of $\\eta$]{$\\eta$-dependence of $\\beta$ factor for cone radii $R=0.4$, $0.7$ and \n$1.0$, measured in the di-jet component of $\\mathtt{jet20}$\nsample.}\\label{eta_correction}\n\\end{center}\n\\end{figure}\n\n\\paragraph{Level-4: Multiple Interactions Corrections}\n\nJet energy measurement is also degraded by the presence of minimum-bias events\nthat come from multiple $p\\bar{p}$ interactions. This correction becomes more \nrelevant at high luminosity, indeed the number of $p\\bar{p}$ interactions is \nPoisson distributed with mean value approximately linear with instantaneous luminosity:\n\\begin{equation} \n\\langle N(\\mathscr{L}\\simeq 10^{32}~\\mathrm{cm}^{-2}\\mathrm{s}^{-1})\\rangle\\simeq3\\mathrm{,}\\quad \\langle N(\\mathscr{L}\\simeq 3\\cdot 10^{32}~\\mathrm{cm}^{-2}\\mathrm{s}^{-1})\\rangle\\simeq 8\\mathrm{.}\n\\end{equation}\nThe energy of particles coming from those processes is estimated from \nminimum-bias events drawing a cone in a random\nposition in the region $0.1<\\eta<0.7$. Figure~\\ref{mp_correction} shows that\nthe measured minimum-bias $E_T$ grows linearly with the number of primary \nvertices\\footnote{Good quality primary vertices are reconstructed through at \nleast 2 COT tracks.}, such quantity, $A_{L4}$, must be subtracted by jet raw \nenergy. The total uncertainty is about $15\\%$, it mostly depends on luminosity \nand event topology.\n\\begin{figure}[t]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{PhysObj\/Figs\/mp_correction}\n\\caption[$E_T$ Correction due to Multiple Interactions]{$E_T$ correction deriving from multiple interactions as a function of primary vertex number for cones with $R=0.4$ (right) and $R=0.7$ (left).}\\label{mp_correction}\n\\end{center}\n\\end{figure}\n\n\\paragraph{Level-5: Absolute Energy Scale Corrections}\n\nWhile $L1$ and $L4$ make jet reconstruction uniform over the whole detector and\nover the global behavior of $p\\bar{p}$ beam interaction, $L5$ correction \n($f_{L5}$) aims to derive, from the detector jet energy measurement, the $E_T$ of particles originating the jet.\n\nThe study is MC driven: first jet events are generated with full CDF detector \nsimulation, then jets are reconstructed both at calorimeter and hadron\ngeneration levels (HEPG) with the use of same clustering algorithm. A \ncalorimeter jet (C) is associated to the corresponding hadron jet (H) if \n$\\Delta R <0.1$. For both HEPG and detector jets the transverse momentum, \n$p_T^C$ and $p_T^H$, is calculated. \nThe absolute jet energy is defined as $\\mathcal{P}(p_T^C|p_T^H)$, the probability\\footnote{Different \n$p_T^H$ can give the same $p_T^C$, in this case the maximum is taken.} to measure $p_T^C$ with a given $p_T^{H}$.\n\nFigure~\\ref{abs_correction} shows the correction factor $f_{L5}$ for different \ncone sizes as function of the different jet transverse energies. The total \nuncertainty is about $3\\%$ and it mainly arises from the determination of\ncalorimetric response to single particles and MC fragmentation modeling.\n\n\\begin{figure}[t]\n\\begin{center}\n\\includegraphics[width=0.85\\textwidth]{PhysObj\/Figs\/abs_correction2}\n\\caption[Absolute Jet Energy Correction as a Function of $p_T$]{Absolute jet energy scale correction and systematic uncertainty ($f_{L5}$) for \\texttt{JETCLU} algorithm with $R=0.4$ cone size.}\\label{abs_correction}\n\\end{center}\n\\end{figure}\n\n\\paragraph{Level-6 \\& Level-7: Underlying Event and Out-of-cone Corrections}\n\nAlthough we do not use L6 and L7 corrections, their description is reported for completeness. They are the last two corrections needed to infer the initial energy of the parton originating the jet. \n\nThe underlying event correction ($L6$) takes into account the interaction \nprocesses which can occur between spectator partons or that originates from initial \nstate radiation (usually soft gluon radiation) while the out-of-cone correction\n($L7$) considers the fraction of particles coming from the original parton that fall outside the jet cone. \n\nThe underlying event energy ($A_{L6}$) must be subtracted to the total jet \nenergy. It was measured studying minimum-bias events during Run I and is \nparametrized with a constant value that scale with the cone radius. Out of cone\nenergy ($A_{L7}$) must be added to the total jet energy, studies are carried \nout with the same jet-parton matching method of $L5$. \n\n\n\n\\section{Neutrino Reconstruction}\\label{sec:metObj}\n\nNeutrinos are the only subatomic particles that leaves the detector completely undetected, therefore their signature is {\\em missing energy}. \n\nSince the longitudinal energies of the colliding partons are unknown and not necessarily equal, we can only say that the total transverse energy of the $p \\bar{p}$ collision is zero. Therefore the total amount of missing transverse energy \\cancel{E}$_T$ (or MET) gives a measurement of the neutrino transverse momentum\\footnote{For a massless neutrino $p_{T}=E_{T}$.} and it is defined as:\n\\begin{equation}\\label{rawmet}\n\\vec{\\mathrm{\\cancel{E}}}_T\\equiv - \\sum_i{\\vec{E}_T^i}\n\\end{equation}\nwhere $\\vec{E}_T^i$ is a vector with magnitude equal to the transverse energy \ncollected by the i-th calorimeter tower and pointing from the interaction\nvertex to the center of the tower. The sum involves all the towers with total \nenergy above $0.1$~GeV in the region $|\\eta|<3.6$.\n\nThe \\mbox{${\\not}{E_T}$} obtained from Equation~\\ref{rawmet}, close to the online reconstructed missing energy, is often referred to as {\\em raw} \\cancel{E}$_T$ (or \\mbox{${\\not}{E_T^{raw}}$}). The fully reconstructed \\mbox{${\\not}{E_T}$} is corrected for the true $z$ vertex position, for the difference between the raw and corrected $E_T$ of the tight and loose jets, for the presence of muons in the event, by subtracting the momenta of minimum ionizing high-$p_T$ muons and adding back the transverse energy of the MIP deposit in the calorimeter towers. These corrections can be summarized in the following equation:\n\\begin{equation}\n\\mbox{${\\not}{E_T}$} = \\mbox{${\\not}{E_T^{raw}}$} - \\sum_{\\text{muon}} p_T + \\sum_{\\text{muon}} (E_T^{EM} + E_T^{HAD}) - \\sum_{jet}(E_T - E_T^{cor}).\n\\end{equation}\n\n\n\n\\section{Secondary Vertex Tagging}\\label{sec:secvtx}\n\nThe algorithms able to select a jet coming from a Heavy-Flavor ($HF$) quark hadronization process are called {\\em $b$-taggers} or {\\em heavy-flavor taggers} and they are of fundamental importance in this and in many other analyses. For example both the top quark and the SM Higgs boson (for $m_H\\lesssim 135$~GeV$\/c^{2}$) have large branching fraction in $b$-quark, therefore an efficient $b$-tagging can dramatically reduce the background of uninteresting physical processes which contain only light-flavor ($LF$) hadrons in their final state.\n\nWe employ the Secondary Vertex Tagger\\footnote{Historically it was the most important \ncomponent in top discovery in 1995.} algorithm (\\verb'SecVtx') to select a $HF$ enriched sample by requiring a $b$-tag on one or both the selected jets. In this way, it is possible to discriminate the $W\\to cs$ and $Z\\to c\\bar{c}\/b\\bar{b}$ decays against the generic $W\/Z\\to u,d, s$ decays. \n\nIn a successive phase of the analysis a flavor separator Neural Network~\\cite{Richter_stop2007,kit_note7816}, developed by the Karsrhue Institute of Technology (also named KIT Flavor Separator), is used to separate jets originating from $c$ and $b$ quarks, allowing a separate measurement of $W\\to cs$ and $Z\\to c\\bar{c}\/b\\bar{b}$.\n\n\n\\subsection{The SecVtx Algorithm}\\label{sec:secvtx_alg}\n\nThe \\verb'SecVtx' algorithm takes advantage of the long life time of $b$-hadrons: a $c\\tau$ value of about $450~\\mu$m together with a \nrelativistic boost due to a momentum of several GeV$\/c$ permits to a $b$-hadron to fly several millimeters\\footnote{The average transverse momentum of a $b$-hadron coming from a $WH$ events is about $40$~GeV$\/c$ for a Higgs boson mass of $120$~GeV$\/c^2$; in that condition a neutral $B^0$ meson of mass $5.28$~GeV$\/c^2$ undergoes a boost $\\beta\\gamma=7.6$ and the average decay length is $3.5$~mm.} away from the primary interaction \nvertex. The relevant quantity is the $c\\tau$ which is approximately the average\nimpact parameter of the outgoing debris of $b$-hadron decays. \nThe decay produces small sub-jets composed by tracks with large impact \nparameter ($d_0$). The silicon detectors (see \nsection~\\ref{track_par}) are able to reconstruct $d_0$ with adequate \nprecision to separate displaced tracks from the prompt tracks coming from the \nprimary interaction. Figure~\\ref{svx_example} shows as a $W+$ jets candidate\nevent with two displaced secondary vertices is identified by \\verb'SecVtx' and \nreconstructed by the CDF event display.\n\\begin{figure}[t]\n\\begin{center}\n\\includegraphics[width=0.7\\textwidth]{PhysObj\/Figs\/svx_example}\n\\caption[$W +$ Jets Candidate Event with Two Reconstructed Secondary Vertices]{$W+$ jets candidate event with two secondary vertices tagged by \n$\\mathtt{SecVtx}$ (run 166063, event 279746). The \\cancel{E}$_T$ direction, a \nmuon track, a prompt track and tracks from the secondary vertices are shown.}\\label{svx_example}\n\\end{center}\n\\end{figure}\n\nTagging is performed for all the jets with $|\\eta|<2.4$ in an event. The algorithm searches for secondary vertices using the tracks within the jet cone of radius \\mbox{$\\Delta R=0.4$} . The usable tracks must satisfy the following requirements:\n\\begin{itemize}\n\\item $p_T>0.5$~GeV$\/c$;\n\\item $|d_0|<0.15$~cm and $|d_0\/\\sigma_0|>2.0$;\n\\item $|z_0-z_{PV}|<2.0$~cm;\n\\item have a minimum number (depending on track reconstruction quality and position) of hits in the silicon detector;\n\\item be seeded or confirmed in the COT;\n\\end{itemize}\na {\\em taggable} jet is defined as a jet containing at least two usable tracks.\n\nThe algorithm works on a two step basis and has two main operation modes\\footnote{An {\\em ultra-tight} operation mode exists but it is rarely used.}, {\\em tight} (the standard one) and {\\em loose}. The operating modes are defined by track and vertex quality criteria~\\cite{secvtx_opt7578} but the two-step selection algorithm remains identical.\n\n\nIn the \\emph{Pass 1} at least three tracks are required to pass loose selection\ncriteria. At least one of the tracks used is required to have $p_T > 1.0$~GeV$\/c$. The selected tracks are \ncombined two by two until a seed secondary vertex is reconstructed, then all \nthe others are added one by one and a quality $\\chi^2$ is computed. Tracks are \nadded or removed depending of their contribute to the $\\chi^2$.\n\nThe \\emph{Pass 2} begins if \\emph{Pass 1} does not find a secondary vertex. Now only \ntwo tracks are required to form a secondary vertex but they must pass tighter requirements: $p_t>1.0$~GeV$\/c$, $|d_0\/\\sigma_0|>3.5$ and one of the tracks must have $p_T > 1.5$~GeV$\/c$.\n\nIf a secondary vertex is identified in a jet, the jet is {\\em tagged}. The bi-dimensional decay length $L_{xy}$ is calculated as the projection on the jet axis, in the $r-\\phi$ \nplane, of the \\verb'SecVtx' vector, i.e. the one pointing from the primary vertex \nto the secondary. The sign of $L_{xy}$ is defined by the angle $\\alpha$ between\nthe jet axis and the \\verb'SecVtx' vector. Figure~\\ref{fig:svx} explains the \ngeometry.\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{PhysObj\/Figs\/svx}\n\\caption[Reconstructed Secondary Vertex: True and Fake $\\mathtt{SecVtx}$]{Schematic representation of $\\mathtt{SecVtx}$ variables and geometry. Left: true reconstructed secondary vertex. Right: negative $\\mathtt{SecVtx}$ tag, falsely reconstructed secondary vertex.}\\label{fig:svx}\n\\end{center}\n\\end{figure}\n\nA secondary vertex coming from a $HF$ hadron is expected to have \nlarge $L_{xy}$. To reduce background due to mismeasured tracks \n\\mbox{$|L_{xy}\/\\sigma_{L_{xy}}|>7.5$} is required\\footnote{Negative $L_{xy}$ has\nno physical meaning but it is important to estimate the mistag probability due \nto resolution effects.}. Other requirements are applied on the \ninvariant mass of the pair of tracks, to avoid $K$ and $\\lambda$ decays, and on\nvertex multiplicity and impact parameter to reject secondary vertices due to \ninteraction with material inside the tracking volume.\n\n\n\\subsection{Tagging Performances and Scale Factors}\\label{sec:btag_sf}\n\nThe performances of a $b$-tagger are evaluated on its efficiency, i.e the rate \nof correctly identified $b$-hadrons over all the produced $b$-hadrons, and on \nits purity, i.e the rate of falsely identified $b$-hadrons in a sample with no \ntrue $b$-hadrons. CDF uses $b\\bar{b}$ QCD MC to evaluate \\verb'SecVtx' efficiency \nrelying on detector and physical processes simulation. Figure~\\ref{fig:tag_eff} \nshows the $b$-tagging efficiency as a function of jet $\\eta$ and $E_T$ for the \ntight and loose \\verb'SecVtx' operating modes. Tagging efficiency drops at large $|\\eta|$ because of tracking acceptance.\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{PhysObj\/Figs\/tag_eff}\n\\caption[\\texttt{SecVtx} Efficiency]{Efficiency of the \\texttt{SecVtx} $b$-tagging algorithm for the tight and loose operating modes.}\\label{fig:tag_eff}\n\\end{center}\n\\end{figure}\n\nAs MC does not reproduce the exact $b$-tagging efficiency of \\verb'SecVtx' a {\\em Scale Factor} ($SF_{Tag}$) is introduced to account for data\/MC difference in the form:\n\\begin{equation}\nSF_{Tag}\\equiv\\frac{\\epsilon_{data}}{\\epsilon_{MC}}.\n\\end{equation}\n\n\nCDF uses two methods to calculate $SF_{Tag}$: both use di-jet samples where one of the two jets contains a low-$p_T$ lepton, electron or muon, which increases the presence of a $HF$ hadron decay, while the other jet can be $b$-tagged by \\verb'SecVtx' or not. \n\nThe method exploiting the electrons~\\cite{eleSF_10178} finds algebraically the relative number of $b$-tagged $HF$-jets with respect to the number of $b$-jets with a semileptonic electron decay. To extract the fraction, it also needs the number of jets with a low-$p_T$ electron but not containing $HF$ hadrons: this is obtained from a jet sample where the $\\gamma\\to e^+e^-$ conversion process is identified, and therefore, with depleted $HF$ content. \n\nThe second method, used as a validation of the first one, exploits the muon semileptonic decay of the $HF$ hadrons~\\cite{muSF_8640}. The algorithm is similar to the electron based one, but the fraction of low-$p_T$ muons not coming from $HF$ hadrons is extracted from a fit of the muon $p_T$ relative to the jet axis (named $p_T^{rel}$) that has a peculiar distribution for $HF$ hadron decays. \n\nThe efficiency on data are therefore compared to the MC results and the $SF_{Tag}$ is extracted.\nFigure~\\ref{fig:tag_sf} shows the $SF_{Tag}$ dependency with respect to the jet $E_T$ as determined by the electron method.\nTable~\\ref{tab:tag_sf} reports the $SF_{Tag}$ for the loose and tight \\verb'SecVtx' operation modes integrated over the variables of interest for the $SF_{Tag}$ parametrization (no strong dependency is seen in any of them).\n The total per-jet tagging efficiency, deconvoluted from tracking effects is about $40$\\% for $b$-jets and $6$\\% for $c$-jets.\n\\begin{table}\\begin{center}\n\\begin{tabular}{cccc}\n\\toprule\n mode & $SF_{Tag}$ & stat. err. & sys. err.\\\\\n\\midrule\nTight & $0.96$ & $0.01$ & $0.05$ \\\\\nLoose & $0.98$ & $0.01$ & $0.05$ \\\\\n\\bottomrule\n\\end{tabular}\n\\caption[\\texttt{SecVtx} Efficiency Scale Factors]{$\\mathtt{SecVtx}$ Scale Factors ($SF_{Tag}$) for the tight and loose $\\mathtt{SecVtx}$ operation modes.}\\label{tab:tag_sf}\n\\end{center}\\end{table}\n\\begin{figure}[!h]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{PhysObj\/Figs\/tag_sf}\n\\caption[\\texttt{SecVtx} Scale Factor]{Evaluation of the \\texttt{SecVtx} $SF_{Tag}$ with the electron method~\\cite{eleSF_10178} and its dependence over the $E_T$ of the jet. The linear fit shows almost no dependency from the $E_T$ of the jet. Table~\\ref{tab:tag_sf} reports the evaluation of $SF_{Tag}$ integrated over the variables of interest for the $HF-$tagger parametrization.}\\label{fig:tag_sf}\n\\end{center}\n\\end{figure}\n\nThe number of falsely \\verb'SecVtx' tagged jets is dubbed \\emph{mistags}.\nMistags can be due to track resolution, long living $LF$ hadrons or \nsecondary interactions with detector material.\n\nThe rate of $W$ plus mistag jets is derived from a sample of events collected with \nan inclusive jet-based trigger with no $HF$ requirement\\footnote{The presence of a \nsmall $HF$ contamination in the sample is a source of systematic uncertainty.}. \nThe mistag parametrization~\\cite{mistag} is obtained from an inclusive jet sample using {\\em negative tags} (see Figure~\\ref{fig:svx}), i.e. $b$-jets which appear to travel back toward the \nprimary vertex. Resolution and material effects are expected to produce false \ntags in a symmetric pattern around the primary interaction vertex. The mistag \nrate is then corrected for the effects of long-lived $LF$ hadrons to take\ninto account the LF contamination giving real secondary vertices. \nA {\\em per-jet} mistag probability, $p^j_{Mistag}$ is parametrized in bins of:\n\\begin{itemize}\n\\item total number of jets in the event, \n\\item jet $E_T$, \n\\item jet $\\eta$, \n\\item track multiplicity within the jet, \n\\item total $E_T$ of the event, \n\\item number of interaction vertices,\n\\item the $z$ vertex position.\n\\end{itemize}\nThis defines the, so called, {\\em Mistag Matrix}. Figure~\\ref{fig:tag_mistag} shows the mistag rate as function of $E_T$ and $\\eta$ of the tagged jets for the tight and loose \\verb'SecVtx'\noperation modes. The approximate per-jet fake rate is about $1$\\% for the tight operation mode and $2$\\% for the loose one.\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{PhysObj\/Figs\/tag_mistag}\n\\caption[$E_T$ and $\\eta$ Mistag Rate]{Rate of wrongly $\\mathtt{SecVtx}$ tagged jets (\\emph{mistags}) as a function of $E_T$ and $\\eta$ of the jets for the tight and loose $\\mathtt{SecVtx}$ operation modes. The rate is derived from an inclusive jet sample.}\\label{fig:tag_mistag}\n\\end{center}\n\\end{figure}\n\n\\section{Neural Network Flavor Separator (KIT-NN)}\\label{sec:kitnn}\n\nThe $b$-tag requirement drastically reduces the contamination of $LF$ jets, however also a per-jet fake rate of $1$\\% can produce enough mistag background to decrease the sensitivity to rare processes. On the other hand a more stringent $b$-tagging requirement would reduce the signal yield. The KIT Flavor Separator Neural Network~\\cite{Richter_stop2007,kit_note7816} (KIT-NN) offers a possible solution to these problems: after a \\verb'SevVtx' tag, the Neural-Network (NN) exploits a broad range of $b$ quark discriminative variables to obtain a continuous distribution with a good separation power between real $b$-jets and $LF$ jets (see Figure~\\ref{fig::kit_nn_out}). \n\n\\begin{figure}[!h]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{PhysObj\/Figs\/kit_yield}\n\\caption[KIT-NN Output Distribution and $b$-jet Efficiency\/Purity Curve]{Left: Output value of the KIT-NN function for $b$, $c$ and $LF$ jets. Right: curve of the $b$-jet efficiency {\\em vs} $b$-jet purity for different selection thresholds of the KIT-NN output.}\\label{fig::kit_nn_out}\n\\end{center}\n\\end{figure}\nBeyond the $LF$ separation power, in this analysis we exploit the KIT-NN distribution in an original way. Thanks to the substantial shape difference between $c$-jets and $b$-jets KIT-NN values, we managed to measure separately $W\\to cs$ and $Z\\to c\\bar{c}\/b\\bar{b}$ contributions (as shown in Chapeter~\\ref{chap:StatRes}).\n\nA short description of the KIT-NN implementation, developed for the single-top search at CDF~\\cite{single_top}, can be useful to understand the physical meaning of such variable. \n\nA Bayesian NN is a supervised learning algorithm~\\cite{BShop_machine_learning} that associates a score to an {\\em event}, according to its likeliness to be {\\em signal} or {\\em background}. An event is defined by an array of $K$ input variables, $\\alpha_i$, and the ouput score, $o \\in [-1,1]$, is obtained from the following function:\n\\begin{equation}\n o =S \\Big( \\sum^H_{j=0} \\omega_j S \\Big( \\sum_{i=0}^{K}\\omega_{i,j} \\alpha_i +\\mu_{0,j} \\Big) \\Big),\n\\end{equation}\nwhere $\\omega_i$, $\\omega_{i,j}$ and $\\mu_{0,j}$ are tunable parameters and $S$ is a sigmoid, or activation, function:\n\\begin{equation}\n S(x) =\\frac{2}{1+e^{-x}} -1,\n\\end{equation}\nThe parameters are derived from a target function optimized over the training set of signal and background labeled events. \n\nThe final result is a per-event signal-or-background {\\em posterior probability} that takes into account non-linear correlations between the input variables. \n\nIn our case, after that a jet has been tagged by \\verb'SevVtx', twenty-five input variables, relative to the secondary vertex and to the jet, are fed into a NeuroBayes\\textsuperscript{\\textregistered} NN~\\cite{neurobayes}. The training signal sample is composed by $t\\bar{t}$, single-top and $W+b\\bar{b}$ MC samples, while the background is built with $W+c\\bar{c}$ and $W+LF$ MC samples. The variables exploit lifetime, mass, and decay multiplicity of the $b$ hadrons using several characteristics of the identified secondary vertex, the properties of the tracks inside the jet and information from the $b$-tag algorithm, for example if {\\em Pass 1} or {\\em Pass 2} reconstruction is used (see Section~\\ref{sec:secvtx_alg}).\nFigure~\\ref{fig::kit_nn_out} shows the obtained classification values for $b$, $c$ and $LF$ jets.\n\nAlthough a careful validation of the input variables was performed, the use of a MC-driven $LF$ training sample introduced a relevant discrepancy in the KIT-NN evaluation of {\\em real} mistagged jets. Therefore a $LF$ correction function was derived from a fake-enriched data sample selected by negative $\\mathtt{SecVtx}$ tags. The final $LF$ KIT-NN distribution is obtained by assigning, to each $LF$-jet, a random value extracted from the corrected $LF$ distribution. The uncorrected template distribution is used as an {\\em optimistic} systematic variation. A similar effect has been hypothized also for $c$-jets but it appears in a less relevan way. For $c$-jets the KIT-NN output is evaluated per-jet each jet and the $LF$-like variation is used as a {\\em pessimistic} systematic variation. Figure~\\ref{fig:kit_sys}\nshows the default (central) KIT-NN distribution and systematic variations for $LF$ and $c$ jets.\n\n\\begin{figure}[!h]\n\\begin{center}\n\\includegraphics[width=0.99\\textwidth]{PhysObj\/Figs\/kit_sys}\n\\caption[KIT-NN Distribution and Systematic Variation for $LF$ and $c$ Jets]{KIT-NN default distribution (continuous lines) and systematic variations (dashed lines) for $LF$-jets (left) and $c$-jets (right).}\\label{fig:kit_sys}\n\\end{center}\n\\end{figure}\n\n\\chapter{Events Selection}\\label{chap:sel}\n\nIn the following we describe the selection requirements applied to produce samples of data enriched in $diboson\\to \\ell\\nu+ HF$ candidates.\n\nEvery physics analysis starts with some kind of candidate signal selection. In the case of collider experiments, it is possible to see three ingredients in it: online (trigger) selection (Section~\\ref{sec:onlineSel}), offline selection (Sections from~\\ref{sec:w_sel} to~\\ref{sec:svm_sel}) and the efficiency evaluation of each of the steps.\n\nThe variety of different physics object that compose the diboson decay channel under investigation is a characteristic of this analysis: charged leptons (electrons, muons and isolated tracks), missing energy, jets and $HF$ jets. These represent a good fraction of all the objects that can be reconstructed at a hadron collider experiment. Furthermore the need to relax the selection requirements in order to increase the statistics and the background-rich $W+2$ jets sample increase the challenges. It worths noticing that the same issues arise also in other primary analyses at CDF: for example single-top or $WH$ searchs.\n\nIn order to deal with these difficulties, one of the major task of this thesis work was the development of a robust analysis framework that could exploit the $\\ell\\nu+HF$ channel in an efficient and reliable way. The package, named $WH$ Analysis Modules (WHAM), is described in Appendix~\\ref{chap:AppWHAM}. Its strengths are modularity, a wide set of configuration options and the use, with improvements, of CDF most advanced tools. The modular and object-oriented approach of the software allowed the handling of the complex selection formed by eleven different lepton types and four data-streams. At the same time it was also possible to develop a set of tools useful also in other contexts: for example a new multivariate multi-jet background rejection tool (see Appendix~\\ref{chap:AppSvm}) or a versatile $b$-tag efficiency evaluation algorithm.\n\nThe actual selection process is divided into two main steps. First we select candidates that have the $\\ell\\nu+2$ jets signature, this defines the {\\em pretag control} sample. Successively we require one or both the jets to be $b$-tagged, obtaining two signal samples ({\\em single} and {\\em double} tagged) enriched in $HF$. Section~\\ref{sec:data_sig} shows the final result of the selection on the CDF dataset and on the diboson signal.\n\n\n\\section{Online Event Selection}\\label{sec:onlineSel}\nThe trigger system, described in Section~\\ref{sec:trigger}, is in charge of the online selection.\n\nIn order to maximize the signal acceptance we use a complex trigger strategy based on several trigger paths collected in four data-streams:\n\\begin{itemize}\n\\item high energy central electrons or \\verb'bhel';\n\\item high energy forward electrons or \\verb'bpel';\n\\item high momentum central muons or \\verb'bhmu';\n\\item high missing transverse energy or \\verb'emet'.\n\\end{itemize}\nAfter online trigger selection but before offline analysis, the quality of the recorded data is crosschecked and the final luminosity is calculated (see Section~\\ref{sec:data-struc}).\n\nA detailed description of the trigger selection, the trigger efficiency ($\\epsilon_{trig}$) evaluation and of the data quality requirements follows in the Section. \n\n\n\\subsection{High Energy Central Electron Trigger}\n\nThe trigger path used for the selection of tight central electron candidate events (CEM) is named \\verb'ELECTRON_CENTRAL_18' and it is stored in the \\verb'bhel' data-stream together with other auxiliary triggers. Single high energy electron identification is based on a calorimeter electromagnetic cluster matched with a reconstructed track.\n\nThe selection proceeds through the three trigger levels, first with minimal requirements and then with more and more sophisticated object reconstruction. L1 requirements are a track with $p_T > 8$~GeV\/c, a central\n($|\\eta| < 1.2$) calorimeter tower with $E_T > 8$~GeV and the ratio between the\nenergy deposited in the hadronic calorimeter to that in the electromagnetic\ncalorimeter ($E^{HAD}\/E^{EM}$) less than $0.125$. At L2, it is required a calorimeter cluster with\n$E_T > 16$~GeV matched to a track of $p_T > 8$~GeV$\/c$. At L3, it requires an\nelectron candidate with $E_T > 18$~GeV matched to a track of $p_T > 9$~GeV$\/c$.\n\nThe \\verb'ELECTRON_CENTRAL_18' trigger path was extensively studied in precision CDF measurements involving $W$ and $Z$ bosons~\\cite{w_z_prl} and now, after each data taking period, the efficiency is evaluated with a standard set of tools~\\cite{perfidia}. A data sample of $W$ boson collected with a \\verb'NO_TRACK' trigger is used to measure the tracking efficiency while a backup trigger is used to measure the efficiency of the calorimeter clustering. The efficiency is \\mbox{$\\epsilon^{CEM}_{trig} =0.982\\pm 0.003$} across the complete dataset with a small $E_T$ and $\\eta$ dependency (see Figure~\\ref{fig:cem_trig}). The trigger turn-on $E_T$ dependency is parametrized with the following function:\n\\begin{equation}\\label{eq:cem_et}\n \\epsilon_{trig}(x) = A - Be^{-c x};\n\\end{equation}\nwhere $x \\equiv E_T$ and $A$, $B$, $c$ are the parameters of the fit. The following equation is used to model the small effect across $\\eta$:\n\\begin{equation}\\label{eq:cem_eta}\n \\epsilon_{trig}(x) = A - \\frac{C}{2\\pi\\sigma} e^{-\\frac{x^2}{2\\sigma^2}};\n\\end{equation}\nwhere $x \\equiv \\eta$ and $A, C, \\sigma$ are the parameters of the fit. Figure~\\ref{fig:cem_trig} shows the parametrization, derived from the full CDF dataset, and applied to MC events.\n\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth, height=0.36\\textwidth]{Sel\/Figs\/cem_et_trig}\n\\includegraphics[width=0.495\\textwidth, height=0.36\\textwidth]{Sel\/Figs\/cem_eta_trig}\n\\caption[\\texttt{ELECTRON\\_CENTRAL\\_18} Trigger Efficiency Parametrization]{\\texttt{ELECTRON\\_CENTRAL\\_18} trigger efficiency parametrization as a function of $E_T$ (left) and $\\eta$ (right) of the candidate electrons. $E_T$ dependence is parametrized by Equation~\\ref{eq:cem_et}, $\\eta$ dependence by Equation~\\ref{eq:cem_eta}. The fit parameters are reported in the legends.}\\label{fig:cem_trig}\n\\end{center}\n\\end{figure}\n\n\n\n\\subsection{High Energy Forward Electron Trigger}\n\nThe candidate high energy forward electrons (in $1.2 < |\\eta_{Det}| < 2.0$ or Plug region) are stored in the \\verb'bpel' data-stream. \n\nElectron triggering in the plug region is particularly challenging because it lacks of the COT coverage and because a large number of soft interactions are boosted in the region closer to the\nbeam-line. A partial solution to these problems is the use of a multiple-object trigger, named \\verb'MET_PEM', dedicated to $W\\to e\\nu$ selection instead of an inclusive electron trigger: there is no track selection but large \\mbox{${\\not}{E_T}$} is required together with a high energy cluster in the forward electromagnetic calorimeter. The L1 requirements are a EM tower of $E_T > 8$~GeV in the forward region and \\mbox{${\\not}{E_T}$} $ > 15$~GeV. L2 confirms \\mbox{${\\not}{E_T}$} requirements but a reconstructed EM cluster of $E_T > 20$~GeV is required and, finally, L3 confirms again the same quantities (\\mbox{\\mbox{${\\not}{E_T}$} $> 15$}~GeV and EM Cluster \\mbox{$E_T > 20$}~GeV) with L3 reconstruction algorithms. \n\nThe trigger shows a turn-on for both the $E_T$ of the cluster and the \\mbox{${\\not}{E_T^{raw}}$}, i.e. corrected only for the position of the primary vertex. The following equation is used for the parametrization of the turn-on:\n\\begin{equation}\\label{eq:phx_trig}\n \\epsilon_{trig}(x) = \\frac{1}{1+ e^{-A(x-B)}};\n\\end{equation}\nwhere $x\\equiv E_T$ or $x\\equiv$\\mbox{${\\not}{E_T^{raw}}$} and $A$ and $B$ are the parameters of the fit, reported in Figure~\\ref{fig:phx_trig}.\n\nThe \\verb'MET_PEM' trigger parametrization was part of my initial work in\nthe CDF collaboration~\\cite{met_pem_trig9493} and is now included in the CDF analysis\ntools~\\cite{perfidia}.\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Sel\/Figs\/phx_et_trig}\n\\includegraphics[width=0.495\\textwidth]{Sel\/Figs\/phx_rawmet_trig}\n\\caption[\\texttt{MET\\_PEM} Trigger Efficiency Parametrization]{\\texttt{MET\\_PEM} trigger efficiency parametrization as a function of $E_T$ (left) of the electron candidate and \\mbox{${\\not}{E_T^{raw}}$} (right) of the event. The \\mbox{${\\not}{E_T^{raw}}$} is corrected only for the position of the primary vertex and therefore it is closer to the trigger level \\mbox{${\\not}{E_T}$} . The $E_T$ and \\mbox{${\\not}{E_T^{raw}}$} dependence is parametrized by Equation~\\ref{eq:phx_trig} and the fit parameters are reported in the legends.}\\label{fig:phx_trig}\n\\end{center}\n\\end{figure}\n\n\n\n\n\\subsection{High Momentum Central Muon Trigger}\n\nCentral tight muon candidate events (CMUP and CMX) are collected by the \\verb'MUON_CMUP18' and the \\verb'MUON_CMX18' triggers and saved into the \\verb'bhmu' data-stream. \n\nFor \\verb'MUON_CMUP18' the main trigger requirements are: at L1, a track (\\mbox{$p_T > 4$}~GeV$\/c$) matched to a stub in both the CMU and CMP detectors. At L2, a confirmed track ($p_T > 15$~GeV$\/c$) and a calorimeter deposit consistent with a MIP in the direction of the track and muon stubs. Finally at L3, a fully reconstructed COT track ($p_T > 18$~GeV$\/c$) whose extrapolation matches the CMU (CMP) hits within $\\Delta x_{CMU} < 10$~cm ($\\Delta x_{CMP} < 20$~cm). \nSimilar requirements are made also for \\verb'MUON_CMX18' but the sub-detector of interest is the CMX (covering $0.6 < |\\eta| < 1.1$) and the L1 requirement is a track with $p_T > 8$~GeV$\/c$ matched to a stub in the CMX chamber. \n\nAs for tight electron candidates, also the central muon triggers were extensively studied~\\cite{w_z_prl} and the efficiency is\nmeasured with a standard procedure~\\cite{perfidia}. Using a data sample of $Z\\to \\mu^+\\mu^-$, a muon passes the trigger requirements and the second muon is checked to pass the trigger selection or not. \n\nAs the CMUP sample represents a relevant fraction of the good quality muon data, an additional effort was done, in the context of the $WH$ analysis~\\cite{hao_trig10747}, to improve the CMUP trigger collection efficiency. Basically we exploited all the CMUP trigger candidates present in the \\verb'bhmu' stream and collected by secondary and auxiliary triggers. The purity of those events was found to be the same of the \\verb'MUON_CMUP18' selected, and we evaluated, with the bootstrapping method, a recover of about $12$\\% efficiency on CMUP online selection. \n\nThe trigger efficiencies are \\mbox{$\\epsilon^{CMUP}_{trig} = 0.976\\pm 0.001$} and \\mbox{$\\epsilon^{CMX}_{trig}=0.873\\pm0.001$}, flat in the kinematic variables. \n\n\n\n\\subsection{High Missing Transverse Energy Triggers}\n\nIn order to recover lepton acceptance lost because of the limited geometrical coverage of the single-lepton triggers, we rely on data collected by the \\texttt{emet} data-stream. \n\nThe \\texttt{emet} stream contains samples with large online \\mbox{${\\not}{E_T}$} or \\mbox{${\\not}{E_T}$} plus reconstructed jets. No charged lepton information is used, therefore the acceptance of $\\ell\\nu + 2$~jets events can be recovered. In particular we use the \\verb'emet' stream to collect Loose Muons (Section~\\ref{sec:LooseMuId}) and Isolated Tracks (Section~\\ref{sec:IsoTrk}) candidates, altogether dubbed Extended Muon Categories (EMC).\n\nThe trigger optimization~\\cite{higgs_ttf8875} was performed with the aim to increase the sensitivity to Higgs boson production, but, as diboson and $WH$ production share the same $HF$ final state, an identical strategy was employed for this analysis. The three main triggers used are: \\verb'MET45', \\verb'MET2J' and \\verb'METDI'. The exact naming and requirements of the trigger paths changed many times along the data taking period due to different instantaneous luminosity conditions but the main classification and characteristics are the following:\n\\begin{itemize}\n\\item \\verb'MET45' requires large \\mbox{${\\not}{E_T}$} in the event, first at L1 and, with increased quality and cut levels, at L2 and L3. The trigger was present along all the data taking with slightly different specifications, however in the last 2\/3 of the data taking the L3 \\mbox{${\\not}{E_T}$} cut was $40$~GeV.\n\\item \\verb'MET2J' requires high \\mbox{${\\not}{E_T}$} in the event at L1 (\\mbox{${\\not}{E_T}$} $ > 28$~GeV), two jets at\nL2 with one of them reconstructed in the central region of the detector ($|\\eta| < 0.9$) and, finally, L3 object confirmation and \\mbox{${\\not}{E_T}$} $ > 35$~GeV. The trigger was Dynamically Pre-Scaled (DPS) to cope with the increasing luminosity in the second half of the data taking.\n\\item \\verb'METDI' requires two jets and \\mbox{${\\not}{E_T}$} in the event, very similarly to \\verb'MET2J', but the L1 requirement, with \\mbox{${\\not}{E_T}$} $ > 28$~GeV and at least one jet, is different.\nThis trigger was introduced after the first $2.4$~fb$^{-\u22121}$ of integrated luminosity.\n\\end{itemize}\n\nThe overlap and composite structure of the \\verb'emet' triggers makes their use less straightforward than the inclusive lepton ones. Furthermore the use of multiple-objects requires to model, on signal and background MC, the possible selection efficiency correlations (e.g. \\mbox{${\\not}{E_T}$} and jets correlations).\nEfficiency studies were performed in~\\cite{Nagai:2010zza, Buzatu:2011zz}. The main results are the definition of appropriate kinematic regions where $\\epsilon_{trig}$ is flat (plateau region) with respect to the jet variables and can be parametrized as a function of the \\mbox{${\\not}{E_T}$} only. \\mbox{${\\not}{E_T}$} is corrected for primary vertex and jet energy but not for the presence of muons. In this way, tight muon candidates, collected from\nother triggers (e.g. \\verb'MUON_CMUP18'), can be used to evaluate the turn-on curves as their contribution to the \\mbox{${\\not}{E_T}$} is not accounted at trigger level. \n\nAn optimal trigger combination strategy for \\texttt{emet} triggers was evaluated in~\\cite{Buzatu:2011zz}. The trigger selection maximizes the $WH$ acceptance after that the \\mbox{$\\epsilon_{trig}(\\mbox{${\\not}{E_T}$} )$} function is evaluated for each data period and kinematic configuration of the three triggers. \n\nThe final combined trigger efficiency reaches $\\epsilon_{trig}\\simeq 0.5$ for $WZ$ simulated events.\n\n\n\n\n\n\\subsection{Data Quality Requirements and Luminosity Estimate}\\label{sec:lumi}\n\nThe last step before offline event selection is the application of strict data quality requirements. These are ensured by the Good Run List (GRL) selection: i.e. it is possible to use only the run periods where the performances of the sub-detectors are optimal and well parametrized.\n\nA conservative approach, used in many of the CDF analyses, is to consider a run as good only if all the sub-detectors used to define the physics objects of the analysis are fully functional. However here, as in the latest $WH$ search results~\\cite{wh94_note10796}, an ad-hoc approach is used~\\cite{ext_grl}. We require a full functionality of the silicon detector (used for jet $b$-tagging) and calorimeters, while, as we are interested in the one-lepton final state, we de-correlate inefficiencies due to shower-max detectors functionality (CES and PES used for CEM and PHX identification) from muon chambers (CMU, CMP, CMX problematic runs). In this way it was possible to recover 8\\% of the collected luminosity.\n\nOnce that the trigger paths and the appropriate GRLs are defined, it is possible to calculate the final integrated luminosity used for the various lepton types:\n\\begin{itemize}\n\\item CEM and PHX electrons use an integrated luminosity of $\\int \\mathcal{L}\\mathrm{dt} = 9.446$~fb$^{\u2212-1}$.\n\\item CMUP muons use an integrated luminosity of $\\int \\mathcal{L}\\mathrm{dt} = 9.494$~fb$^{-\u22121}$. \n\\item CMX muons use an integrated luminosity of $\\int \\mathcal{L} \\mathrm{dt}= 9.396$~fb$^{-1}$.\n\\item EMC charged leptons use an integrated luminosity of $\\int \\mathcal{L}\\mathrm{dt} = 9.288$~fb$^{-\u22121}$.\n\\end{itemize}\nThe calculated luminosity has systematic uncertainty of $6$\\% (see Section~\\ref{clc_par}).\n\n\n\\section{Offline Event Selection}\n\nNow, on good data, we apply the object identification algorithms described in Chapter~\\ref{chap:objects} to select the desired final state of $\\ell\\nu + HF$ jets. \n\nTable~\\ref{tab:cuts} summarizes the complete selection {\\em cut flow} and the following Sections (from \\ref{sec:w_sel} to~\\ref{sec:svm_sel}) describe in detail the selection criteria. Part of the cut flow is common to many other CDF analyses while the applied multi-jet rejection algorithm, as well as some other specific choices, are completely original.\n\nBecause we rely on simulation to evaluate the signal acceptance, the data\/MC agreement is cross checked for each one of the employed algorithms with data driven methods. Well known physics processes are selected in control samples and the efficiency on data ($\\epsilon^{Data}_{alg}$) and MC ($\\epsilon^{MC}_{alg}$) is estimated. The ratio between them, with appropriate correlations and systematic errors evaluation, is a correction Scale Factor ($SF_{alg}$) to be applied on the simulation:\n\\begin{equation}\n SF_{alg} =\\frac{ \\epsilon^{MC}_{alg}}{\\epsilon^{Data}_{alg}}\n\\end{equation}\nThe $SF_{alg}$ associated to the lepton identification ($SF_{l}$) and to the $b$-tagging ($SF_{Tag}$) are particularly important in this analysis because of the large number of different lepton reconstruction algorithms and because of the $HF$ final state of the signal. The $SF$s are discussed along with the selection criteria in the following Sections\n\n\n\\begin{table}[h]\n \\begin{center}\n \\begin{tabular}{c}\n \\toprule\n {\\bf Selection Cuts}\\\\\\midrule\n\n Trigger Fired \\\\\n Good Run Requirement\\\\\n $|z_0^{PV}| \\le 60$~cm \\\\\n Charged Lepton Selection \\\\\n $|z_0^{lep} \u2212- z_0^{PV}| < 5$~cm\\\\\n Di-Lepton Veto\\\\\n $Z\\to \\ell\\ell$ Veto\\\\\n Cosmic Ray Veto \\\\\n $\\gamma\\to e^+e^-$ Conversion Veto\\\\\n \\mbox{${\\not}{E_T}$} $ > 15$~GeV\\\\\n 2 Tight Jets Selection \\\\\n $M_{Inv}(jet1, jet2)>20$~GeV$\/c^2$\\\\\n 1 or 2 \\texttt{SecVtx} Tags\\\\\n {\\bf Multi-jet Rejection:} \\\\\n {\\em SVM}$_{CEM, EMC}>0$, {\\em SVM}$_{CMUP, CMX}>-0.5$, {\\em SVM}$_{PHX}>1$\\\\\n \\bottomrule\n \\end{tabular}\n \\caption[Selection Summary]{Summary of all the selection requirements applied in the analysis.}\\label{tab:cuts}\n \\end{center}\n\\end{table}\n\n\n\n\n\\subsection{$W\\to\\ell\\nu$ Offline Selection}\\label{sec:w_sel}\n\nThe leptonic decay of the $W$ boson is selected by requiring exactly one high energy charged lepton and missing transverse energy signaling the presence of the neutrino.\n\nThe correct application of these two basic requirements exploits several of the identification algorithms described in the previous Chapter and other selection criteria in order to purify the lepton sample.\nThe summary of all the selection steps, applied on data and MC, is the following:\n\\begin{itemize}\n\\item identification of the primary interaction vertex (Section~\\ref{sec:primVtx}) within the interaction fiducial region of 60 cm from the detector center ($|z_0^{PV}| \\le 60$~cm). Minimum bias events are used to evaluate on data the efficiency of this cut. The result, averaged over all the data taking periods, is\\footnote{The initial number of MC events, used to normalize all the simulated acceptances, is derived after this cut so that $\\epsilon_{z_0}$ should be applied to obtain the final yields.}:\n \\begin{equation}\n \\epsilon_{z_0} = 0.9712\\pm 0.0006.\n \\end{equation}\n\\item Identification of one charged and isolated lepton candidate originating\nfrom the primary vertex ($|z_0^{lep} \u2212- z_0^{PV}| < 5$~cm). A charged lepton candidate is one of the eleven lepton identification algorithms described in Sections from~\\ref{sec:eleSel} to~\\ref{sec:IsoTrk}: CEM, PHX (tight electrons), CMUP, CMX (tight muons), BMU, CMU, CMP, CMIO, SCMIO, CMXNT (loose muons) and ISOTRK (isolated tracks). This requirement does not produce any appreciable signal loss.\n\\item Veto of di-leptonic candidate events: if an event contains two tight or loose electrons or muons both isolated or not, it is rejected as a Drell-Yan event or $t\\bar{t}$ dileptonic candidate. Events with an ISOTRK candidate are vetoed if reconstructed together with another ISOTRK or a tight lepton. An efficiency of approximately 90\\% is estimated from signal MC.\n\\item Further rejection of $Z\\to\\ell\\ell$ candidates is obtained for CEM, PHX, CMUP and CMX leptons by vetoing events in which the tight lepton can be paired with a track or a EM calorimeter cluster that form an invariant mass within the range $[76, 116]$~GeV$\/c^2$. An efficiency of approximately 90\\% is estimated from signal MC.\n\\item A cosmic tagger was implemented within the CDF offline code~\\cite{cosmic_tag}\nto reject high-$p_T$ cosmic muons which interact with the detector simultaneously with a bunch crossing. The algorithm exploits muon chambers, COT and TOF, calorimeter energy and timing information, providing\nalmost 100\\% rejection of cosmic ray events, with a negligible loss of signal efficiency.\n\\item Rejection of electrons originating from photon conversions, $\\gamma\\to e^+e^{-\u2212}$,\nis obtained by applying the CDF conversion tagger~\\cite{conversion_tag}. The algorithm looks for two opposite-sign tracks (one of them belonging to the identified electron) and requires $\\Delta \\cot(\\theta) \\le 0.02$ and the distance at the closest approach between them $D_{xy} \\le 0.1$~ cm. The rejection efficiency of the conversion veto is about $65$\\% with a $1\\div 2$\\% signal loss.\n\\item Finally, \\mbox{${\\not}{E_T}$} $ > 15$~GeV is required to signal the presence of a neutrino. The \\mbox{${\\not}{E_T}$} is fully corrected for position of the primary interaction vertex, presence of jets and muons in the event (see Section~\\ref{sec:metObj}). Approximatly 7\\% of the signal is lost with this requirement.\n\\end{itemize}\n\nFor each run period and lepton identification algorithm, the agreement\nbetween data and MC is measured comparing the $Z\\to\\ell\\ell$ data sample against the simulation. Table~\\ref{tab:lepId} reports, averaged on all the run range, the yield correction, $SF_{l}$ that covers the\ndata\/MC differences. The main differences arise in the simulation of the isolation and of the muon chamber response.\n\n \\begin{table}[h]\n \\begin{center}\n \\begin{tabular}{cccc}\n \\toprule\n\n CEM & PHX & CMUP & CMX \\\\\n $0.973 \\pm 0.005$ & $0.908 \\pm 0.009$ & $0.868\\pm 0.008$ & $ 0.940\\pm 0.009$\\\\\n \\end{tabular}\n \n \\begin{tabular}{cccccc}\n \\midrule\n BMU & CMU & CMP & SCMIO & CMIO & ISOTRK \\\\\n\n $ 1.06\\pm 0.02$ & $ 0.88\\pm 0.02$ & $ 0.86\\pm0.01$ & $ 1.02\\pm 0.01$ & $ 0.97 \\pm 0.02$ & $ 0.94 \\pm 0.04$ \\\\\n \\bottomrule \n \\end{tabular}\n\n \\caption[Lepton Algorithm Data\/MC Scale Factors ($SF_l$)]{Data\/MC Scale Factors ($SF_l$) for each of the single lepton identification algorithms used in the analysis, tight charged leptons in the first row and EMC categories in the lower row.}\\label{tab:lepId}\n \\end{center}\n \\end{table}\n\nIf several lepton identification algorithms cover the same detector region they can not be considered independent and a further correction\\footnote{This correction was introduced for the first time in this analysis and in the $WH\\to \\ell\\nu+ b\\bar{b}$ search~\\cite{wh94_note10796}.} to the simulation is needed. As shown in Figure~\\ref{fig:split_alg}, loose muon and ISOTRK largely overlap with the tight lepton selection and, often, a MC lepton candidate is reconstructed by more than one lepton selection algorithm. Because of this and to avoid ambiguity, we apply a prioritized selection, i.e. we label a lepton candidate with the highest priority algorithm according to:\n \\begin{equation}\n\\begin{gathered}\n CEM > PHX > CMUP > CMX > BMU > CMU >\\\\\n CMP > SCMIO > CMIO > CMXNT > ISOTRK\n\\end{gathered}\n \\end{equation} \n\n\nIn most of the cases, selection cuts are orthogonal and only one lepton reconstruction is possible, however, if this is not true {\\em and} $SF_{l}\\ne 1$, the same simulated lepton has a probability greater than zero of not being identified by the highest priority algorithm\\footnote{Quality criteria and historical reasons are at the base of the prioritization of the algorithms. Prioritization is indirecly present also in the $SF_{l}$ measurements: for example in the ISOTRK $SF_{l}$ measurement the lepton candidates are checked not to be reconstructed as tight leptons, while the check is not performed in the opposite case, i.e. no ISOTRK identification is performed during the tight $SF_{l}$ measurement.}. In this case the MC event under consideration is classified under more lepton categories with different weights, according to:\n\\begin{equation}\\label{eq:composite_sf}\n \\omega_{l} = SF_{l} + SF_{l}\\prod_{l'>l} \\big(1- SF_{l'} \\big)\n\\end{equation}\nwhere $\\omega_{l}$ is a composite $SF_{l}$. Equation~\\ref{eq:composite_sf} reduces to $\\omega_{l} = SF_{l}$ when a lepton candidate is reconstructed only by one identification algorithm, otherwise it takes into account the probability of {\\em not} identification from algorithms with higher priority ($l'>l$).\n\nFigure~\\ref{fig:eta_multisf} shows the corrected $\\eta$ distribution for the EMC lepton category (combination of loose muon and ISOTRK algorithms) while Figure~\\ref{fig:eta_pre_post} shows the complete background estimate distribution for the combined $\\eta$ of the EMC category before and after the correction of Equation~\\ref{eq:composite_sf}.\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.8\\textwidth]{Sel\/Figs\/emc_eta_sf}\n\\caption[$\\eta$ Distribution of the $\\omega_{l}$ Correction on EMC]{Effect of the $\\omega_{l}$ correction described in Equation~\\ref{eq:composite_sf} across the $\\eta$ distribution of the EMC lepton category selected on $WZ$ MC. The correction appears as a factor $f_{\\omega_l}\\ne 1$ for each event where the reconstructed lepton is identified by more than one algorithm.}\\label{fig:eta_multisf}\n\\end{center}\n\\end{figure}\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Sel\/Figs\/WZ_Pretag_ISOTRK_LepEta_noSFCorr}\n\\includegraphics[width=0.495\\textwidth]{Sel\/Figs\/WZ_Pretag_ISOTRK_LepEta_wSFCorr}\n\\caption[Effect of $\\omega_{l}$ on EMC $\\eta$ Distribution after Background Estimate]{$\\eta$ distribution of the EMC lepton category with full background evaluation before (left) and after (right) the application of the multiple $SF_{l}$ correction defined by Equation~\\ref{eq:composite_sf}.}\\label{fig:eta_pre_post}\n\\end{center}\n\\end{figure}\n\n\n\n\n\n\\subsection{Offline Jets Selection}\\label{sec:jetSel}\n\nThe final aim of this analysis is the study of the invariant mass spectrum of a two jet, $HF$ enriched, final state. \n\nJet selection is based on the \\verb'JETCLU' $R = 0.4$ algorithm, described in Section~\\ref{sec:jetObj}, and the \\texttt{SecVtx} $HF$-tagger algorithm, described in Section~\\ref{sec:secvtx}.\nOne of the goals of this analysis is to demonstrate the ability of the CDF experiment to identify a resonance in the background rich sample of $W+HF$ jets. Therefore the jet selection was kept simple with the aim to understand the sample itself. The specific requirements are the following:\n\\begin{itemize}\n\\item selection of exactly two tight jets. Jets central in the detector and with large $E_T$ (after L5 energy correction) are named tight jets if:\n\n\\begin{equation}\n E^{jet,cor}_T > 20\\mathrm{~GeV}\\quad\\textrm{and}\\quad |\\eta^{jet}_{Det}| < 2.0;\n\\end{equation}\n\n\\item a minimal lower bound on the invariant mass distribution of the jet pair: $M_{Inv}(jet1, jet2) > 20$~GeV$\/c^2$;\n\\item one or both of the selected jets $HF$-tagged by the \\texttt{SecVtx} algorithm;\n\\item jets are ordered according to their $E_T^{jet,cor}$ and classified as: $jet1$ and $jet2$. \n\\end{itemize}\nNo other requirements are applied on the jets and, in particular, we do not apply any criteria to the loose jets present in the event. \n\nThe $HF$-tag requirement divides the $\\ell\\nu+ 2$ jets sample in three analysis regions:\n\\begin{description}\n\\item[Pretag:] no $HF$-tag requirement is applied and, due to the small $HF$ contamination, the sample is used as a control region (see Section~\\ref{sec:wjets_pretag}).\n\\item[Single-tag:] events with exactly one jet tagged as $HF$ compose this signal region. The statistics of the selected sample is quite large therefore it is important for the sensitivity of this analysis to $WW\\to\\ell\\nu + c\\bar{s}$ decay.\n\\item[Double-tag:] events with both jets tagged as $HF$. The statistics of this sample is small because of the strict double-tag requirement, however a large fraction of the sensitivity for the $WZ\\to\\ell\\nu + b\\bar{b}$ signal comes from this region as the background contamination is also low.\n\\end{description}\n\nAccording to Section~\\ref{sec:btag_sf}, the $HF$ selection efficiency of the \\texttt{SecVtx} algorithm must be corrected with an appropriate data\/MC $SF_{Tag}$. Furthermore a contamination due to $LF$ quarks producing secondary vertices will contribute to the signal sample. The $HF$ and $LF$ components for $1$\u2212-Tag and $2$--Tag selection are accounted in the signal MC with a per-event tagging probability, $\\omega_{k-\u2212Tag}$ obtained by the combination of per-jet tagging probabilities:\n\\begin{equation}\\label{eq:1tag}\n \\omega_{1-Tag}=\\sum_{i=1}^{n}p^i_{Tag}\\cdot\\Big(\\prod^n_{j=1,j\\ne i}\\big( 1 - p^j_{Tag}\\big)\\Big),\n\\end{equation}\n\\begin{equation}\\label{eq:2tag}\n \\omega_{2-Tag}=\\sum_{i=1}^{n-1}p^i_{Tag}\\cdot\\Big( \\sum_{j>i}^{n}p^i_{Tag} \\cdot \\Big( \\prod^n_{h=1,h\\ne i,h\\ne j}\\big( 1 - p^k_{Tag}\\big)\\Big)\\Big),\n\\end{equation}\nwhere the per-jet tagging probability is defined by $p^j_{Tag} \\equiv SF_{Tag}$, if the jet is matched\\footnote{The matching is satisfied if $\\Delta R(jet, had) < 0.4$ where we test the fully reconstructed jet 4-vector against all the final state hadrons 4-vectors.} to a $HF$ hadron, or $p^j_{Tag} \\equiv p^j_{Mistag}$, if the event is matched to a $LF$ hadron.\n\nEquations~\\ref{eq:1tag} and~\\ref{eq:2tag} are applied to obtain \\texttt{SecVtx} 1-Tag or 2-Tags probabilities. A further generalization of the previous equations was implemented in the WHAM analysis framework (see Appendix~\\ref{chap:AppWHAM}) so that efficiencies can be calculated for any $b$-tagging algorithm (and combinations of them), any tag and jet multiplicity.\n\n\n\\subsection{Multi-Jet Background Rejection}\\label{sec:svm_sel}\n\nMulti-jet events can fake the $W$ boson signature when one jet passes the high $p_T$ lepton selection criteria and fake \\mbox{${\\not}{E_T}$} is generated through jets mis-measurement. Although the probability for such an event is small, the high rate of multi-jet events, combined with the small cross section of the processes of interest, make this an important background. Additional difficulties arise in the simulation of the mixture of physics and detector processes contributing to this background. Therefore it is desirable to reduce the multi-jet events as much as possible.\n\nFor this purpose, an original and efficient multi-variate method based on the Support Vector Machine (SVM) algorithm was developed for this analysis.\n\nThe basic concept behind a SVM classifier is quite simple: given two sets of $n$-dimensional vectors (the training sets) the algorithm finds the {\\em best separating hyper-plane} between the two. The plane is then used again to classify a newly presented vector in one of the two sets. Non linear separation can be achieved with the use of appropriate transformations called {\\em Kernels}. SVM classification performs better than other multivariate algorithms, like Neural Networks, in the case of statistically limited training sets and shows more stability with respect to over-training~\\cite{svmbook}.\n\nThe work described here improves a previously developed software package\\footnote{The previously created SVM based multi-jet rejection~\\cite{CHEP_svm} was already tested and employed in previous analyses~\\cite{wh75_note10596, wh94_note10796, cdfDiblvHF_Mjj2011}.}, based on the \\texttt{LIBSVM}~\\cite{libsvm} library, able to perform algorithm training, variable ranking, signal discrimination and robustness test. Here we report only a summary of the results while a detailed description of optimization, training procedure and variable definition is given in Appendix~\\ref{chap:AppSvm}. \n\nWe trained two specific SVMs, one for the {\\em central} and one for the {\\em forward} region of the detector. \nThe SVMs aimed to improve the purity of the $W\\to e\\nu+2$ jets sample, however we avoided the use of input variables related to the electron identification (i.e. cluster or track related) while we focused on the kinematic variables. This gave an additional result as the multi-jet rejection proved to be optimal also for other lepton identification algorithms. For example the central SVM was used on all the lepton identification algorithms in the region $|\\eta|<1.2$. \n\nAppendix A describes in details all the input variables used in the trainings.\nThe SVM for the {\\em central} region exploits the following eight input variables:\n\\begin{itemize}\n\\item $W$ related variables: $M^W_T$, \\mbox{${\\not}{E_T^{raw}}$}, ${\\not}{p_T}$, $\\Delta \\phi(lep,$ \\mbox{${\\not}{E_T}$} $)$, $\\Delta R(\\nu^{min},$ $lep)$;\n\\item global variables: $MetSig$, $\\Delta \\phi(jet1,$ \\mbox{${\\not}{E_T}$} ), $\\Delta \\phi({\\not}{p_T},$ \\mbox{${\\not}{E_T}$} ).\n\\end{itemize} \nThe specific SVM for the {\\em forward} region of the detector (used only for the PHX electron selection) exploits six input variables:\n\\begin{itemize}\n\\item $W$ related variables: $M^W_T$, \\mbox{${\\not}{E_T^{raw}}$}, ${\\not}{p_T}$; \n\\item global variables: $MetSig$, $\\Delta \\phi({\\not}{p_T},$ \\mbox{${\\not}{E_T^{raw}}$}), $\\Delta \\phi({\\not}{p_T},$ \\mbox{${\\not}{E_T}$} ).\n\\end{itemize}\n\nTable~\\ref{tab:cuts} shows different SVM selection values for the different lepton types:\n\\begin{equation}\n\\textrm{{\\em SVM}}_{CEM, EMC}>0\\textrm{;}\n\\end{equation}\n\\begin{equation}\n \\textrm{{\\em SVM}}_{CMUP, CMX}>-0.5\\textrm{;}\n\\end{equation}\n\\begin{equation}\n \\textrm{{\\em SVM}}_{PHX}>1\\textrm{.}\n\\end{equation}\nWe choose different thresholds due to the different probabilities to fake a $W$ signature: lower for the tight muon selection (CMUP, CMX), higher for CEM electrons and EMC lepton category, and way higher for PHX electrons, identified in a detector region with low track reconstruction efficiency and large calorimeter occupancies. \n\nFigure~\\ref{fig:svm_templates} shows the SVM output distribution for CEM and PHX leptons for the signal and background MC samples as well as for the multi-jet background models (see Section~\\ref{sec:qcd} for their description).\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.495\\textwidth]{Sel\/Figs\/svm_cem} \\includegraphics[width=0.495\\textwidth]{Sel\/Figs\/svm_phx}\n\\caption[SVM Output Distribution for CEM and PHX Leptons]{SVM output distribution for CEM (left) and PHX (right) leptons for the signal and background MC samples as well as for the multi-jet background model (see Section~\\ref{sec:qcd} for the description).}\\label{fig:svm_templates}\n\\end{center}\n\\end{figure}\n\nTable~\\ref{tab:svm_eff} summarizes the signal selection efficiencies for the SVM cuts in the different lepton categories; all the analysis selection cuts are applied except for the tagging requirement. The evaluation of the final multi-jet background contamination is described in Section~\\ref{sec:qcd} and a summary of the results is reported in Tables~\\ref{tab:wjet_pretag_qcd_fit} and~\\ref{tab:wjet_tag_qcd_fit}.\n\nThe training procedure requires the SVM to pass several quality criteria about data\/MC agreement. We also checked the stability of the selection efficiency. We found it stable against a variation of up to 5\\% of the SVM input variables. Therefore we did not find necessary to add a specific SF for the SVM selection.\n\n \\begin{table}\n\\begin{center}\n \\begin{tabular}{cccccc}\n \\toprule\n Lepton & CEM & PHX & CMUP & CMX & EMC \\\\\\midrule\n $\\epsilon_{SVM}^{WW}$ & $0.95 \\pm 0.02$ & $0.82 \\pm 0.03$ & $0.97 \\pm 0.03$ & $0.97 \\pm 0.04$ & $0.94 \\pm 0.03$ \\\\\n $\\epsilon_{SVM}^{WZ}$ & $0.96 \\pm 0.03$ & $0.83 \\pm 0.04$ & $0.98 \\pm 0.04$ & $0.99 \\pm 0.05$ & $0.95 \\pm 0.03$ \\\\\n \\bottomrule\n \\end{tabular}\n \\caption[$WW$ and $WZ$ SVM Selection Efficiency]{$WW$ and $WZ$ signal selection efficiencies for the SVM cuts used in the lepton selection. Reported errors are statistical only.}\\label{tab:svm_eff}\n\\end{center}\n \\end{table}\n\n\n\n\\section{Final Signal Estimate and Data Selection}\\label{sec:data_sig}\n\nThe complete event selection can now be applied to both the full data sample and the signal MC. \n\nThe final signal yield is derived from the following equation, for each signal process ($Sig=WW$, $WZ$, $ZZ$), each lepton identification algorithm ($l$) and for single and double tagged ($k=1$, $2$) selections:\n\\begin{equation}\\label{eq:mc_estimate}\n N^{Sig} = \\epsilon_{l, k}^{Sig}\\cdot\\sigma_{p\\bar{p}\\to Sig}\\cdot\\int\\mathcal{L}_{l}\\mathrm{dt};\n\\end{equation}\nthe factors are:\n\\begin{itemize}\n\\item the cross sections, $\\sigma_{p\\bar{p}\\to Sig}$, for inclusive diboson production, reported in Table~\\ref{tab:mc_dib}, are evaluated at NLO~\\cite{mcfm_prd,mcfm_url};\n \\begin{table}\n \\begin{center}\n \\begin{tabular}{ccc}\n \\toprule\n Signal & $\\sigma$ (pb), NLO & Initial Events \\\\\n \\midrule\n $WW$ & 11.34 $\\pm$ 0.66 & 4.6 M\\\\ \n $WZ$ & 3.47 $\\pm$ 0.21 & 4.6 M\\\\ \n $ZZ$ & 3.62 $\\pm$ 0.22 & 4.8 M\\\\ \n \\bottomrule\n \\end{tabular}\n \\caption[Diboson MC Information and NLO Cross Section]{MC information about the diboson signal. The reported cross sections are evaluated at NLO precision~\\cite{mcfm_prd,mcfm_url}. \\texttt{PYTHIA}~\\cite{pythia} v6.216 generator and parton-showering is used for the MC simulation.}\\label{tab:mc_dib}\n \\end{center}\n \\end{table}\n\n\\item the integrated luminosity, $\\int\\mathcal{L}_{l}dt$, is calculated in Section~\\ref{sec:lumi} for the different leptons\/data-streams;\n\\item $\\epsilon_{l, k}^{Sig}$ is the simulated selection acceptance accounting for trigger and $|z_0|$ cut efficiencies ($\\epsilon_{trig}$, $\\epsilon_{z_0}$), lepton and $b$-tag efficiency corrections ($\\omega_l$, $\\omega_{k-Tag}$):\n \\begin{equation}\\label{eq:all_w_tag}\n \\epsilon_{l, k}^{Sig} = \\epsilon_{trig}\\cdot\\epsilon_{z_0}\\cdot\\omega_{l}\\cdot\\omega_{k-Tag}\\cdot\\frac{N^{Sig}_{l,k}}{N^{Sig}_{z_0}}\n \\end{equation}\nwhere $N^{Sig}_{l,k}$ is the number of MC events passing the lepton and tag selection and $N^{Sig}_{z_0}$ is the number of events generated in the fiduciality region of the detector. Equation~\\ref{eq:all_w_tag} simplifies slightly in case we are evaluating the pretag region yield:\n\\begin{equation}\\label{eq:all_w_pretag}\n \\epsilon_{l, pretag}^{Sig} = \\epsilon_{trig}\\cdot\\epsilon_{z_0}\\cdot\\omega_{l}\\cdot\\frac{N^{Sig}_{l,pretag}}{N^{Sig}_{z_0}}\n\\end{equation}\nwhere $N^{Sig}_{l,pretag}$ is the number of MC events passing the lepton selection at pretag level. \n\\end{itemize}\nThe final result of the event selection is reported in Table~\\ref{tab:data_sig} for the CDF dataset and the signal expectation; also the pretag control region is shown for comparison. \nThe selection suppresses the $Z\\to\\ell\\ell$ events making the total $ZZ$ contribution almost negligible ($\\approx 3$\\% of the total diboson yield), however this component is included in the signal because no distinction is possible between the $Z\\to HF$ decay originating from $WZ$ or $ZZ$ production.\n\n\\begin{table}[h]\n \\begin{small}\n \\begin{center}\n \\begin{tabular}{cccccc}\n \\toprule\n \\multicolumn{6}{c}{\\bf Pretag Selection}\\\\\n Lepton & CEM & PHX & CMUP & CMX & EMC \\\\\\midrule\n {\\bf Data} & {\\bf 80263} &{\\bf 27759} & {\\bf 39045} & {\\bf 22465} & {\\bf 35810} \\\\ \n $WW$ & 2012.4 $\\pm$ 168.7 & 705.5 $\\pm$ 59.3 & 1047.7 $\\pm$ 88.5 & 537.0 $\\pm$ 45.5 &1157.7 $\\pm$ 115.3 \\\\\n $ZZ$ & 22.4 $\\pm$ 1.9 & 6.1 $\\pm$ 0.5 & 25.2 $\\pm$ 2.2 & 13.2 $\\pm$ 1.1 &37.5 $\\pm$ 3.8\\\\\n $WZ$ & 307.0 $\\pm$ 26.2 & 135.7 $\\pm$ 11.6 & 168.6 $\\pm$ 14.5 & 95.3 $\\pm$ 8.2 &225.7 $\\pm$ 22.8\\\\\n \\midrule\n \\multicolumn{6}{c}{\\bf Single-Tag Selection}\\\\\n Lepton & CEM & PHX & CMUP & CMX & EMC \\\\ \\midrule\n {\\bf Data} & {\\bf 3115} &{\\bf 1073 } & {\\bf 1577} & {\\bf 830} & {\\bf 1705} \\\\ \n\n $WW$ & 84.35 $\\pm$ 11.8 & 25.05 $\\pm$ 3.54 & 43.7 $\\pm$ 6.17 & 23.68 $\\pm$ 3.35 & 53.11 $\\pm$ 8 \\\\\n $ZZ$ & 1.85 $\\pm$ 0.19 & 0.21 $\\pm$ 0.02 & 2.45 $\\pm$ 0.25 & 1.37 $\\pm$ 0.14 & 3.45 $\\pm$ 0.39 \\\\ \n $WZ$ & 29.2 $\\pm$ 2.95 & 12.32 $\\pm$ 1.22 & 16 $\\pm$ 1.66 & 9.21 $\\pm$ 0.94 & 20.54 $\\pm$ 2.39 \\\\\n\n\\midrule \n \\multicolumn{6}{c}{\\bf Double-Tag Selection}\\\\\n Lepton & CEM & PHX & CMUP & CMX & EMC \\\\\\midrule\n {\\bf Data} &{\\bf 175 } & {\\bf 62 } & {\\bf 92 } & {\\bf 49 } & {\\bf 126} \\\\ \n $WW$ & 0.72 $\\pm$ 0.19 & 0.18 $\\pm$ 0.05 & 0.35 $\\pm$ 0.09 & 0.2 $\\pm$ 0.05 & 0.49 $\\pm$ 0.13 \\\\\n $ZZ$ & 0.26 $\\pm$ 0.04 & 0.03 $\\pm$ 0.01 & 0.46 $\\pm$ 0.06 & 0.29 $\\pm$ 0.04 & 0.63 $\\pm$ 0.10 \\\\ \n $WZ$ & 5.28 $\\pm$ 0.75 & 2.6 $\\pm$ 0.37 & 2.52 $\\pm$ 0.36 & 1.67 $\\pm$ 0.24 & 3.52 $\\pm$ 0.54 \\\\\n \\bottomrule\n\n \\end{tabular}\n \\caption[Final Signal Estimate and Data Selection]{Final result of the event selection for the $W+2$ jets dataset and signal expectation in the different lepton categories. The pretag control region is shown together with single and double-tagged signal regions. The uncertainties from MC statistics, productionn cross section, luminosity, lepton identification, trigger and $HF$-tagging are considered in the table.}\\label{tab:data_sig}\n \\end{center}\n \\end{small}\n \\end{table}\n\n\\chapter{Statistical Analysis and Results}\\label{chap:StatRes}\n\nThe measurement of a potential signal (and its properties) over a predicted background requires a statistical analysis of the selected events.\n\nThe detection of diboson events in the $\\ell\\nu+HF$ final state is a challenging problem as it combines a small expected signal yield, a sizable irreducible background and large systematic uncertainties typical of the hadronic environment.\n\nAfter the full event selection (Chapter~\\ref{chap:sel}), the remaining background processes are mainly of irreducible nature ($W+c\/c\\bar{c}\/b\\bar{b}$) and still overwhelm the signal by more than a factor twenty. \n\nAs the diboson production is a resonant process while $W+HF$ is not, a feasible and optimal signal extraction strategy is the shape analysis of the di-jet invariant mass distribution, $M_{Inv}(jet1,jet2)$.\n\nThe use of the di-jet mass has two effects. Firstly, we can not distinguish $WZ\\to\\ell\\nu+b\\bar{b}\/c\\bar{c}$ decay from $ZZ\\to\\ell\\not{\\ell}+b\\bar{b}\/c\\bar{c}$ decay where one lepton is lost. Therefore, from now on, the $WZ$ diboson production is considered together with the small amount of selected $ZZ$ events\\footnote{Table~\\ref{tab:data_sig} shows that $ZZ$ is about $1\/10$ of the $WZ$ expected events and 1\/30 of the total diboson sample.}. A more important effect is that, due to the low mass resolution, typical of the hadronic final states, we can not separate contribution from $WW\\to \\ell\\nu+c\\bar{s}$ from the $WZ\\to \\ell\\nu+b\\bar{b}\/c\\bar{c}$, in the single-tagged event selection\\footnote{In the double-tagged event selection the only signal contribution comes from $WZ\\to \\ell\\nu+b\\bar{b}$.}. We solve this problem with the use of a second variable, the NN flavor-separator (KIT-NN) described in Section~\\ref{sec:kitnn}. We use it on single-tagged events to build a bi-dimensional distribution:\n\\begin{equation}\n M_{Inv}(jet1,jet2)\\textrm{ {\\em vs }KIT-NN} \\textrm{.}\n\\end{equation}\nThis, together with the simple $M_{Inv}(jet1,jet2)$ distribution of the double-tagged events, allows to measure $WW$ and $WZ\/ZZ$ contributions separately.\n\nThe statistical analysis is performed with a methodology and a software tool, named {\\em mclimit}~\\cite{mclimit, single_top}, also used in the CDF and Tevatron Higgs searches~\\cite{cdf_cmb, tev_cmb} and adapted to the present analysis.\n\nThe technique, described in Section~\\ref{sec:likelihood}, is based on the integration, over data, of a binned likelihood function where the diboson production cross section is a free parameter. The systematic uncertainties, listed in Section~\\ref{sec:sys_desc}, are included with a Bayesian approach and marginalized to increase the sensitivity. \n\nThe likelihood is built by dividing the selected data in eight orthogonal channels depending on the kinematic properties and on the background composition:\n\\begin{itemize}\n\\item {\\bf 4 Lepton categories} are derived from the data-streams used to collect the events (see Section~\\ref{sec:onlineSel}): central electrons (CEM), forward electrons (PHX), central muons (CMUP+CMX) and extended muons (EMC). The kinematic is homogeneous within each sample.\n\\item {\\bf 2 Tag categories} are derived from the single or double \\texttt{SecVtx} tagged event selection for each lepton category. The sample composition drastically changes in the two regions. Section~\\ref{sec:mjj_kit_distrib} shows the $M_{Inv}(jet1,jet2)$ distribution in the double-tagged categories and slices of the bi-dimensional distribution, $M_{Inv}(jet1,jet2)$ {\\em vs} KIT-NN, used for the single-tagged categories. Figure~\\ref{fig:2d_templates} shows four examples (for the CEM lepton selection) of the bi-dimensional templates used for the single-tagged channels evaluation: $WW$, $WZ\/ZZ$ signals and $W+c\\bar{c}$, $W+b\\bar{b}$ backgrounds.\n\\end{itemize}\n\\begin{figure}\n \\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/ww_templ}\n \\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/wzzz_templ}\\\\\n \\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/wcc_templ}\n \\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/wbb_templ}\n \\caption[Example of $M_{Inv}(jet1,jet2)$ {\\em vs } KIT-NN Templates]{Four examples of the bi-dimensional templates, $M_{Inv}(jet1,jet2)$ {\\em vs } KIT-NN, used for the statistical analysis of the single-tagged channels in the CEM lepton selection: $WW$ signal (top left), $WZ\/ZZ$ signals (top right), $W+c\\bar{c}$ background (bottom left), $W+b\\bar{b}$ background (bottom right). The contribution of $c$ and $b$ quarks cluster at opposite sides of the KIT-NN distribution, the diboson production cluster around an invariant mass resonant peak while the generic $HF$ production spreads along the $M_{Inv}(jet1,jet2)$ direction.}\\label{fig:2d_templates}\n\\end{figure}\n\nThe measured diboson cross sections, $\\sigma_{WW}$ and $\\sigma_{WZ\/ZZ}$, are obtained in Section~\\ref{sec:cx_measure} from the evaluation of the Bayesian posterior distribution of the combined likelihood. \n\nAs a first result, we measure the total diboson production cross section, fixing the $WW$ and the $WZ\/ZZ$ relative contributions to the SM prediction. Then we remove the constraint on the $WW$ and the $WZ\/ZZ$ relative fractions. We let them free to float independently so to obtain a measurement where the full correlation between $\\sigma_{WW}$ and $\\sigma_{WZ\/ZZ}$ is accounted. Finally we also obtain two individual measurements integrating out $\\sigma_{WW}$ or $\\sigma_{WZ\/ZZ}$ one at the time from the Bayesian posterior distribution.\n\nIn Section~\\ref{sec:sigma_eval}, we discuss the significance of the different measurements and we show the evidence of the total diboson production. The previous version of this analysis~\\cite{cdfDiblvHF_Mjj2011}, performed with a dataset of $7.5$~fb$^{-1}$ and reported in Appendix~\\ref{App:7.5}, was the first evidence of this process at a hadron collider.\n\nIn general the result confirms the SM prediction for diboson production also when they are detected in the $HF$ final state. This provides additional confidence on the capability and the methods used by the CDF collaboration for the Higgs search in the $H\\to b\\bar{b}$ channel. \n\n\n\\section{Construction of the Likelihood Function}\\label{sec:likelihood}\n\nThe comparison of observed data with the templates derived from signal and background is possible with the use of a likelihood function in which we use a Bayesian approach to include the prior probabilities of the background and systematic effects.\n\nWe assume that the outcome of the data observation $n_{i,j}$, in the $j^{th}$ bin of the $i^{th}$ input histogram follows the Poisson statistics and the expectation value, $\\mu_{i,j}$, depends by the estimated backgrounds and signals. Therefore the complete likelihood function, $\\mathscr{L}$, of the bin by bin outcome is:\n\\begin{equation}\\label{eq:likelihood}\n \\mathscr{L}(\\vec{\\alpha}, \\vec{s}, \\vec{b}| \\vec{n}, \\vec{\\beta}) =\\prod^{Hists}_{i} \\prod^{Bins}_{j} \\frac{e^{-\\mu_{i,j}}\\mu_{i,j}^{n_{i,j}}}{n_{i,j}!}\n\\end{equation}\nwhere the first product runs over the total number of input channels, $Hist$ (i.e. the different histograms provided for them), and the second runs on the bins of each distribution. The bin expectation value contains the signal and background dependence: \n\\begin{equation}\\label{eq:mu_likelihood}\n \\mu_{i,j}(\\vec{\\alpha}, \\vec{s}, \\vec{b}| \\vec{\\beta}) = \\sum_p^{Sgn} \\alpha_{p} s_{i,j,p}(\\vec{\\beta}) + \\sum_q^{Bkg} b_{i,j,q}(\\vec{\\beta})\\textrm{,}\n\\end{equation}\nwhere $s_{i,j,p}$ represents the $p^{th}$ unknown signal, $b_{i,j,q}$ is the $q^{th}$ estimated background process and $\\alpha_{p}$ is the {\\em scaling parameter} used to measure the amount of the different signals. The unknown expectation of the signal is parametrized by flat, uniform, positive prior distribution in $\\vec{\\alpha}$. The last element of Equations~\\ref{eq:likelihood} and~\\ref{eq:mu_likelihood} is $\\vec{\\beta}$, the vector of the nuisance parameters: it incorporates in the likelihood the systematic uncertainties. They are parametrized as fractional variations on the $s_{i,j,p}$ and $b_{i,j,q}$ rates: \n\\begin{equation}\n s_{i,j,p}(\\vec{\\beta})= s_{i,j,p}^{central}\\prod_{k}^{Sys}(1+\\frac{\\sigma_{i,j,k}}{s_{i,j,p}^{central}}\\beta_k)\n\\end{equation}\n\\begin{equation}\n b_{i,j,p}(\\vec{\\beta})= b_{i,j,p}^{central}\\prod_{k}^{Sys}(1+\\frac{\\sigma_{i,j,k}}{b_{i,j,p}^{central}}\\beta_k) \n\\end{equation}\nwhere $\\frac{\\sigma_{i,j,k}}{s_{i,j,p}^{central}}$, and $\\frac{\\sigma_{i,j,k}}{b_{i,j,p}^{central}}$ are the relative uncertainties related to the systematic effect $k$ and the $\\beta_k$ variables. It is important to notice that we distinguish:\n\\begin{itemize}\n\\item {\\em total rate uncertainties}, like the luminosity dependence of the signal, where there is no bin by bin variation of the systematic effect: $\\sigma_{i,j,k}\\equiv \\sigma_{i,k}$.\n\\item {\\em shape uncertainties} where the rate variation can change on a bin by bin basis, evaluated by the ratios between the central and the varied histograms.\n\\end{itemize}\nTo account for the nuisance parameters effect and correlate them across different channels we introduce in Equation~\\ref{eq:likelihood}, for each systematic, a Gaussian probability constraint centered in zero and with unitary variance. The final result is:\n\\begin{equation}\\label{eq:likelihood_full}\n \\mathscr{L}(\\vec{\\alpha}, \\vec{s}, \\vec{b}| \\vec{n}, \\vec{\\beta}) =\\prod^{Hists}_{i} \\prod^{Bins}_{j} \\frac{e^{-\\mu_{i,j}}\\mu_{i,j}^{n_{i,j}}}{n_{i,j}!}\\prod_{k}^{Sys}e^{-\\beta_k^{2}\/2}\\textrm{,}\n\\end{equation}\nthe previous equation is the full likelihood used in the analysis of the results.\n\n\\subsection{Parameter Measurement and Likelihood Integration}\\label{sec:likelihood_meas}\n\nBy maximizing the likelihood, we obtain the {\\em measurement} of the unknown $\\vec{\\alpha}$.\n\nThe maximum of Equation~\\ref{eq:likelihood_full} can be found in different ways: with a fit in the multidimensional space of $\\vec{\\alpha}$ and $\\vec{\\beta}$ or with the integration ({\\em marginalization}) of the nuisance parameters over their prior probability distributions. \n\nThe second technique, exploited here for the measurements, returns a {\\em Bayesian posterior distribution} of $\\vec{\\alpha}$: the minimal extension that covers $68$\\% of the distribution around the maximum gives the one-standard-deviation confidence band and the uncertainty on the measurement.\n\nSection~\\ref{sec:cx_measure} reports the results of the marginalization, where the integration of the nuisance parameters is performed numerically with a Markov-chain adaptive integration~\\cite{mk-chain_book}.\n\nHowever we perform also a fit of the likelihood function because it gives the {\\em best} outcome for all the nuisance parameters. Those values, the result of the fit and the reduced $\\chi^2$, should be consistent with the results obtained by the marginalization.\n\n\n\n\\section{Sources of Systematic Uncertainties}\\label{sec:sys_desc}\n\nAlong the preceding Chapters, we discussed several sources of rate and shape systematics and we need to include them in the likelihood Equation~\\ref{eq:likelihood_full}. \n\nThe {\\em rate-only} systematic effects that we include are summarized in the following points:\n\\begin{description}\n %\n\\item{\\em Initial and Final State Radiation} uncertainties are estimated by changing (halving and doubling) the parameters related to ISR and FSR emission on the signal MC. Half of difference between the two shifted samples is taken as the systematic uncertainty on the signal samples. The total effect on the signal acceptance is about $4\\%$.\n\\item{\\em Parton Distribution Functions} uncertainties are evaluated by reweighting each event of the signal MC according to several PDFs parametrisation and to the generator level information of the event. Then the acceptance is evaluated again giving a rate variation with respect to the original PDFs. The exact procedure is described in~\\cite{pdf_note7051}. The effect on the signal acceptance was evaluated in previous similar analyses to be around $1\\div 2$\\%. It is added in quadrature with the ISR\/FSR systematic as they both influence only the signal acceptance.\n\\item{\\em $b$-tag Scale Factor} uncertainty comes from the measured \\texttt{SecVtx} $b$-tag efficiency variation. The $SF_{b-Tag}$ is $0.96\\pm0.5$ for a single $b$-matched and \\texttt{SecVtx} tagged jet. The uncertainty, propagated through the per-event tagging probability, is assumed to be double in the case of $c$-jets. The uncertainty is applied on the signal and on all the EWK backgrounds described in Section~\\ref{sec:mc_bkg}.\n\\item{\\em Luminosity} measurement uncertainty contributes for an overall 6\\% rate uncertainty on the signal and on all the EWK backgrounds.\n\\item{ \\em Lepton Acceptance} uncertainty derives from the quadrature sum of trigger efficiency measurements and lepton ID Scale Factors. They range from $1 \\div 2$\\% for tight lepton categories (CEM, CMUP+CMX, PHX) to $6$\\% for the EMC leptons collected by the \\mbox{${\\not}{E_T}$} plus jets triggers. These uncertainties are applied on signal and EWK samples.\n\\item {\\em Top Production} uncertainty is a $10$\\% rate uncertainty applied on all the top-related processes ($t\\bar{t}$, single-top $s$ and $t$ channels). It covers the cross section theoretical prediction uncertainty and the acceptance differences due to systematic top quark mass variation. \n\\item {\\em QCD} normalization, independent for each lepton channel, is constrained to be within 40\\% of the value extracted in Section~\\ref{sec:qcd}.\n\\item{\\em Mistag} uncertainty is derived from the mistag matrix $\\pm 1\\sigma$ variation and propagated to the final $W+LF$ sample. The rate variations are $11$\\% and 21\\% for single and double-tagged channels respectively.\n\\item{\\em K-Factor} uncertainty is of 30\\% on $W+HF$ rate estimate. It also includes the $b$-tag SF uncertainty used to derive this correction. We consider an uncorrelated uncertainty for $W+c\\bar{c}\/W+b\\bar{b}$ and $W+c$ backgrounds as the first process is produced by strong interaction while the latter is of electroweak nature.\n\\item {\\em $Z+$Jets} normalization uncertainty is, conservatively, set to $45$\\% as it includes the uncertainty on $Z+HF$ production.\n\\end{description}\nAppendix~\\ref{app:rate_sys} summarizes the systematic rate variation on the different templates.\n\nWe also account for {\\em shape} systematic variations from three sources: JES, $Q^2$ and KIT modeling.\n\\begin{itemize}\n\\item {\\em Jet Energy Scale} shape uncertainty (JES) is estimated by shifting the JES of the input templates by $\\pm 1\\sigma$ from the nominal value. The acceptance of the process is allowed to change, therefore a new background estimate is performed with the JES varied templates. This produces a simultaneous rate and shape uncertainty. All the templates are affected except the fake-$W$ sample that is derived from data.\n\\item {\\em $W+$ Jets $Q^2$ Scale} uncertainty is obtained by halving and doubling the nominal generation $Q^{2}$ (defined in Equation~\\ref{eq:scale}) of all the $W+$ jets samples ($W+b\\bar{b}$, $W+c\\bar{c}$, $W+c$, $W+LF$). Again, we have a shape and rate uncertainty as a new background estimate is performed with the varied templates.\n\\item {\\em KIT-NN} shape uncertainty is the last systematic that we take into account. We apply two kind of independent variations: one on the fake-$W$ templates and the other on $c$-quarks and mistag templates. \nFor the multi-jet template, a variation of the flavor composition is applied allowing more or less $b$-quark like events in the fake-$W$ KIT-NN template. For templates with relevant $c$-quark component ($WW$, $WZ$, $ZZ$, $W+c$ and $W+c\\bar{c}$) and $W+LF$ a correction to the KIT-NN output is derived from negative tags in data. This systematic variation influences only the shape and not the rate of the final discriminant because it leaves untouched the background composition.\n\n\\end{itemize}\n\nIn the case of templates with few events passing all requirements\\footnote{The MC events can be spread across a large number of bins, especially in the bi-dimensional case where each template is composed by $50\\times 4$ bins.} large statistical fluctuations can introduce a bias in the evaluation of the shape systematics. Therefore, shape variations are filtered to reduce the statistical noise. Filters are widely employed in modern image processing for different purposes. We choose a {\\em median filter smoothing}~\\cite{median_filter} because it maintains long range correlations among the histogram bins. We use a 5-bin filter along the $M_{Inv}(jet1,jet2)$ direction for the bi-dimensional, single-tagged, templates while a 3-bin filter is used for double-tagged templates\\footnote{A double bin width is used for double-tagged $M_{Inv}(jet1,jet2)$ templates w.r.t. single-tagged ones as a 5-bin long range correlation would spoil the systematic variation.}. We also apply a low and high boundary to the possible template variations to reduce them in $[0.5,2.0]$ range.\nFigure~\\ref{fig:smoothing} shows an example of the double-tagged CEM lepton $WZ$ template before and after the smoothing filtering.\n\n\n\\begin{figure} [!h]\n\\centering\n\\includegraphics[width=0.4\\textwidth,height=0.35\\textwidth]{StatRes\/Figs\/noSmooth_SVTSVT_CEM}\\hspace{0.7cm}\n\\includegraphics[width=0.4\\textwidth,height=0.35\\textwidth]{StatRes\/Figs\/smooth_SVTSVT_CEM}\n\\caption[Example of $M_{Inv}(jet1,jet2)$ Template Before and After Smoothing]{Example of the double-tagged CEM lepton $WZ$ $M_{Inv}(jet1,jet2)$ template rate variation due to JES systematic before (left) and after (right) a 3-bin median filter smoothing~\\cite{median_filter}. The red (blue) line indicates the one sigma down (up) JES variation.}\\label{fig:smoothing}\n\\end{figure}\n\n\n\n\n\n\\section{$M_{Inv}(jet1,jet2)$ and KIT-NN Distributions}\\label{sec:mjj_kit_distrib}\n\nThe final templates, including all the systematic effects, are then used to build the likelihood of Equation~\\ref{eq:likelihood_full}. \n\nWe combine a total of eight different channels: {\\em four} lepton sub-samples (CEM, PHX, Tight Muons, EMC) times {\\em two} $b$-tag prescriptions (single and double \\texttt{SecVtx} tags). For the double \\texttt{SecVtx} tagged events, the signal discrimination is based only on the di-jet invariant mass, $M_{Inv}(jet1,jet2)$, while for the single-tagged events we exploit the bi-dimensional distribution of $M_{Inv}(jet1,jet2)$ {\\em vs} KIT-NN flavor separator. The KIT-NN output ranges from $-1$ to $1$ and it is divided in four equal size bins: the rightmost is highly enriched in $b$-like jets, while the others have variable composition of $b$-like, $c$-like and $LF$-like jets. \n\nAs noted in Section~\\ref{sec:likelihood_meas}, we first obtain the maximum value of the likelihood with a fit. The best fit result allows to check the agreement of the predicted distribution with data. The reduced $\\chi^2$ of the fit is:\n \\begin{equation}\n \\frac{\\chi^2}{ NDoF}=\\frac{744.4}{664} = 1.12\\textrm{,}\n \\end{equation}\ncorresponding to a probability $P=0.984$. To further investigate the agreement of data and prediction, several post-fit distributions are shown from Figure~\\ref{fig:mjj_2tag_nokit} to~\\ref{fig:mjj_1tag_allkit}. We evaluate on them both the $\\chi^2$ the Kolmogorov-Smirnov tests.\n\\begin{itemize}\n\\item Figure~\\ref{fig:mjj_2tag_nokit} shows the $M_{Inv}(jet1,jet2)$ distribution for double-tagged events;\n\\item Figure~\\ref{fig:kit_1tag} shows the KIT-NN distribution for single-tagged events after the integration of all the $M_{Inv}(jet1,jet2)$ values used in the bi-dimensional distribution;\n\\item Figure~\\ref{fig:mjj_1tag_allkit} shows the $M_{Inv}(jet1,jet2)$ distribution for the single-tagged channel, added for all the lepton categories and integrated across all the KIT-NN values as well as the two most interesting KIT-NN regions: KIT-NN$>0.5$, $b$-enriched, and KIT-NN$<0.5$ with contribution from $c$ and $LF$ quarks.\n\\end{itemize}\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.99\\textwidth]{StatRes\/Figs\/AddSelected_DT_TIG_EMC_PHX}\n\\caption[$M_{Inv}(jet1,jet2)$ Distribution for 2 \\texttt{SecVtx} Tag Events]{$M_{Inv}(jet1,jet2)$ distribution for candidate events with two \\texttt{SecVtx} tags, all the lepton categories have been added together. The best fit values for the rate and shape of the backgrounds are used in the figure. } \n\\label{fig:mjj_2tag_nokit}\n\\end{figure}\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.99\\textwidth]{StatRes\/Figs\/AddSelected_ST_MjjKitProj_TIG_EMC_PHX}\n\\caption[KIT-NN Distribution for 1 \\texttt{SecVtx} Tag Events]{KIT-NN distribution for candidate events with a single \\texttt{SecVtx} tag, all the lepton categories have been added together. The best fit values for the rate and shape of the backgrounds are used in the figure.} \n\\label{fig:kit_1tag}\n\\end{figure}\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.99\\textwidth]{StatRes\/Figs\/AddSelected_ST_KitAll_TIG_EMC_PHX}\\\\\\vspace{0.5cm}\n\\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/AddSelected_ST_KitLow_TIG_EMC_PHX}\n\\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/AddSelected_ST_KitHigh_TIG_EMC_PHX}\n\\caption[$M_{Inv}(jet1,jet2)$ Distribution for 1 \\texttt{SecVtx} Tag Events]{$M_{Inv}(jet1,jet2)$ distribution for candidate events with a single \\texttt{SecVtx} tag, all the lepton category have been added together. Top distribution shows the results for all KIT-NN values while the bottom distributions are separated for KIT-NN$<0.5$ (bottom left) and KIT-NN$>0.5$ (bottom right). The best fit values for the rate and shape of the backgrounds are used in the figures.}\n\\label{fig:mjj_1tag_allkit}\n\\end{figure}\n\n\n\n\n\n\n\\section{Cross Section Measurement}\\label{sec:cx_measure}\n\nThe actual cross section measurement is performed by marginalizing the likelihood with respect to the nuisance parameters and studying the resulting Bayesian posterior as a function of the diboson signal cross section.\n\nFirst we measure the total diboson cross section, $\\sigma_{Diboson}^{Obs}$, constraining the relative $WW$ and $WZ\/ZZ$ cross sections to the SM prediction. The resulting Bayesian posterior distribution is shown in Figure~\\ref{fig:posterior_dib} together with the 68\\% and 95\\% confidence intervals. The measured cross section is:\n\\begin{equation}\n \\sigma_{Diboson}^{Obs}=\\left(0.79\\pm 0.28\\right)\\times \\sigma_{Diboson}^{SM} = \\left(14.6\\pm 5.2\\right)\\textrm{~pb,}\n\\end{equation}\nwhere the errors include statistical and systematic uncertainties and $\\sigma^{SM}_{Diboson}$ is the SM predicted cross section derived from Table~\\ref{tab:mc_dib}:\n\\begin{equation}\n \\sigma_{Diboson}^{SM}= \\sigma^{SM}_{WW}+ \\sigma^{SM}_{WZ\/ZZ} = \\left( 18.43\\pm0.73\\right)\\textrm{~pb.}\n\\end{equation}\n\nIn order to separate the different components, we exploit the $c$ versus $b$ classification power of KIT-NN and the different sample composition of single and double-tagged events to obtain a separate measurement of $WW$ and $WZ\/ZZ$. \n\nWe iterate the cross section measurement procedure but, this time, $\\sigma_{WW}$ and $\\sigma_{WZ\/ZZ}$ are left free to float independently (i.e. not constrained to the SM ratio). Figure~\\ref{fig:posterior_2d} shows the resulting Bayesian posterior distribution. The maximum of the posterior distribution gives the value of the measured cross sections:\n\\begin{equation}\n\\sigma^{Obs,2D}_{WW}=\\left(0.50^{+0.51}_{-0.46}\\right)\\times\\sigma^{SM}_{WW} = \\left(5.7^{+5.8}_{-5.2}\\right)\\textrm{~pb}\n\\end{equation}\nand\n\\begin{equation}\n \\sigma^{Obs,2D}_{WZ\/ZZ}=\\left(1.56^{+1.22}_{-0.73}\\right)\\times\\sigma^{SM}_{WZ\/ZZ}= \\left(11.1^{+8.7}_{-6.3}\\right)\\textrm{~pb,}\n\\end{equation}\nwhere the SM predictions are $\\sigma^{SM}_{WW} = 11.34\\pm0.66$~pb, $\\sigma^{SM}_{WZ\/ZZ} =7.09\\pm0.30$~pb and the errors are evaluated by the intersection of the $x$ and $y$ position of the maximum with the boundary of the smallest area enclosing the 68\\% of the posterior distribution. The smallest areas enclosing 68\\%, 95\\% and 99\\% of the posterior integrals give the contours of one, two and three standard deviations and are explicitly shown in Figure~\\ref{fig:posterior_2d} with correlation between $\\sigma^{Obs,2D}_{WW}$ and $\\sigma^{Obs,2D}_{WZ\/ZZ}$.\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.9\\textwidth]{StatRes\/Figs\/1d_posterior_blessSyle}\n\\caption[Bayesian Posterior Distribution of $\\sigma_{Diboson}$ Measurement]{The Bayesian posterior, marginalized over nuisance parameters, is shown. The maximum value is the central value of the cross-section, the blue and azure areas represent the smallest intervals enclosing 68\\%, 95\\% of the posterior integrals, respectively. The final cross section measurement is\n\\mbox{$\\sigma_{Diboson}^{Obs}=\\left(0.79\\pm 0.28\\right)\\times \\sigma_{Diboson}^{SM} = \\left(14.6\\pm 5.2\\right)\\textrm{~pb}$}.}\n\\label{fig:posterior_dib}\n\\end{figure}\n\nThe $WW$ and $WZ\/ZZ$ channels are also analyzed separately by projecting the two-dimensional Bayesian posterior on the $\\sigma_{WW}$ and the $\\sigma_{WZ\/ZZ}$ axes. In this way, the two processes are considered as background one at the time. For both $WW$ and $WZ\/ZZ$ we re-computed the maximum values and confidence intervals. Figure~\\ref{fig:posterior_2d_proj} shows the results, the measured cross sections are:\n\\begin{equation}\n\\sigma^{Obs}_{WW}=\\left(0.45^{ +0.35}_{ -0.32}\\right)\\times\\sigma^{SM}_{WW} = \\left(5.1^{+4.0}_{-3.6}\\right)\\textrm{~pb}\n\\end{equation}\nand\n\\begin{equation}\n \\quad \\sigma^{Obs}_{WZ\/ZZ}=\\left(1.64^{+0.83}_{-0.78}\\right)\\times\\sigma^{SM}_{WZ\/ZZ} = \\left(11.6^{+5.9}_{-5.5}\\right)\\mathrm{~pb,} \n\\end{equation}\nwhere the errors include statistical and systematic uncertainties.\n\n\\begin{figure}[h!]\n \\centering\n\\includegraphics[width=0.99\\textwidth]{StatRes\/Figs\/2d_posterior}\n\\caption[2-Dim Bayesian Posterior Distribution of $\\sigma_{WW}$ and $\\sigma_{WZ\/ZZ}$ Combined Measurement]{The Bayesian posterior, marginalized over nuisance parameters (scaled to SM expectation), is shown in the plane $\\sigma_{WW}$ {\\em vs} $\\sigma_{WZ\/ZZ}$. The measured cross sections correspond to the maximum value of $\\sigma^{Obs,2D}_{WW}=0.50\\times\\sigma^{SM}_{WW} = 5.7\\textrm{~pb}$ and $\\sigma^{Obs,2D}_{WZ\/ZZ}=1.56\\times\\sigma^{SM}_{WZ\/ZZ}= 11.1\\textrm{~pb}$. The red, blue and azure areas represent smallest areas enclosing 68\\%, 95\\% and 99\\% of the posterior integrals, respectively.}\n\\label{fig:posterior_2d}\n\\end{figure}\n\n\\begin{figure}[h!]\n\\centering\n\\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/1d_WW_posterior_blessSyle}\n\\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/1d_WZZZ_posterior_blessSyle}\n\\caption[Bayesian Posterior Distributions of $\\sigma_{WW}$ and $\\sigma_{WZ\/ZZ}$ Separate Measurements]{The Bayesian posterior, function of $\\sigma_{WW}$ and $\\sigma_{WZ\/ZZ}$ marginalized over nuisance parameters, is shown after projection on the $\\sigma_{WW}$ (left) and $\\sigma_{WZ\/ZZ}$ (right) axes.\nThe maximum value is the central value of the cross-section, the blue and azure areas represent the smallest intervals enclosing 68\\%, 95\\% of the posterior integrals, respectively. The final cross section measurements are: \\mbox{$\\sigma^{Obs}_{WW}=\\left(0.45^{ +0.35}_{ -0.32}\\right)\\times\\sigma^{SM}_{WW} = \\left(5.1^{+4.0}_{-3.6}\\right)\\textrm{~pb}$} and \\mbox{$\\quad \\sigma^{Obs}_{WZ\/ZZ}=\\left(1.64^{+0.83}_{-0.78}\\right)\\times\\sigma^{SM}_{WZ\/ZZ} = \\left(11.6^{+5.9}_{-5.5}\\right)\\mathrm{~pb}$}.}\\label{fig:posterior_2d_proj}\n\\end{figure}\n\n\n\\section{Evaluation of the Statistical Significance}\\label{sec:sigma_eval}\n\nTo compute the significance of the measurements we perform a {\\em hypothesis test} comparing data observation to the {\\em null hypothesis} ($H_0$).\n\nRandom generated Pseudo Experiments (PEs) are extracted from the predicted background processes distribution in the $H_0$ hypothesis (i.e. excluding the diboson production): this is straightforward once we know the probability distribution of the background and of the nuisance parameters. Then we repeat the cross section measurements with the complete marginalization of the likelihood. We expect a distribution peaking at $\\sigma_{Diboson}^{PEs}\/\\sigma_{Diboson}^{SM} =0$ and we compare it to the measured cross section.\n\nFigure~\\ref{fig:significance_dib} shows the possible outcomes of many cross section PEs in a back\\-ground-only and in a background-plus-signal hypothesis. The number of times that a background fluctuation produces a cross section measurement greater than \\mbox{$\\sigma^{Obs}_{Diboson}=0.79\\times\\sigma^{SM}_{Diboson}$} has a $p-$value of $0.00209$. \n\nThe result is an evidence for the diboson production in $\\ell\\nu +HF$ final state with a significance\\footnote{A two sided significance estimate is used because both upper and lower fluctuation of the cross section measurement are considered in the integral of the {\\em null} hypothesis cross section distribution above the measured value.} of $3.08\\sigma$.\n\nThen, we evaluate the single $WW$ and $WZ\/ZZ$ significances in a similar way: PEs are generated with null hypothesis for both $WW$ and $WZ\/ZZ$ signals. Then, the cross section PEs measurements are projected along the $\\sigma_{WW}$ vs $\\sigma_{WZ\/ZZ}$ axes and compared with $\\sigma^{Obs}_{WZ}$ and $\\sigma^{Obs}_{WZ\/ZZ}$. The result of the $p-$value estimates are shown in Figure~\\ref{fig:significance_2d_proj}. \n\nWe obtain: $p-$value$_{WW}=0.074565$ and $p-$value$_{WZ\/ZZ}=0.011145$. They correspond to a significance of $1.78\\sigma$ and $2.54\\sigma$ for $WW$ and $WZ\/ZZ$ respectively.\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.9\\textwidth]{StatRes\/Figs\/significance_diboson}\n\\caption[Diboson Significance Evaluation]{Possible outcomes of many diboson cross section measurements from Pseudo Experiments (PEs) generated in a background-only and in a background-plus-signal hypothesis. The $p-$value for $\\sigma^{Obs}_{Diboson}=0.79\\times\\sigma^{SM}_{Diboson}$ is $0.00209$, corresponding to a significance of $3.08\\sigma$.}\n\\label{fig:significance_dib}\n\\end{figure}\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/significance_WW}\n\\includegraphics[width=0.495\\textwidth]{StatRes\/Figs\/significance_WZZZ}\n\\caption[$WW$ and $WZ$ Separate Significance Evaluation]{Possible outcomes of many diboson cross section measurements from Pseudo Experiments (PEs) generated in a background-only and in a background-plus-signal hypothesis in the $\\sigma_{WW}$ vs $\\sigma_{WZ\/ZZ}$ plane and then projected on the $\\sigma_{WW}$ (left) and $\\sigma_{WZ\/ZZ}$ (axis). The $p-$values of $0.074565$ and $0.011145$ correspond to a significance of $1.78\\sigma$ and $2.54\\sigma$ for $WW$ and $WZ\/ZZ$ respectively.}\n\\label{fig:significance_2d_proj}\n\\end{figure}\n\n\n\n\\chapter{Theoretical and Experimental Overview}\\label{chap:physics}\n\n\nThe goal of particle physics is the understanding of the principles of Nature.\n\nThe quest is pursued through the scientific method: the observation of a phenomenon is explained by a hypothesis that must be, successively, verified or rejected by experimental evidences.\n\nIn this prospect, the observed phenomenon is the existence itself of the atomic and sub-atomic structure of matter, the hypothesis is the Standard Model theory of Elementary Particles and Fundamental Interaction (SM) while the experimental tools are the high energy physics colliders and detectors, available nowadays. \n\nThe main infrastructure of the SM~\\cite{pdg2010, glasgow, weinberg_lep} was developed in the 70's and, since then, it showed to be a very successful theory. One of its main success was the prediction of new elementary particles, later observed at hadron collider experiments. The discovery of the $W$ and $Z$ force carrier vector bosons~\\cite{ua1_w_obs, ua2_w_obs, ua1_z_obs,ua2_z_obs} and of the $top$ quark~\\cite{CDF_top_obs,d0_top_obs} shed light on the fundamental structure of the matter. \n\nStrengthened by these results, SM describes the {\\em electromagnetic}, {\\em weak} and {\\em strong} interactions, three of the four fundamental forces that compose the physics description of Nature. The fourth force, gravitational interaction, is left out but it is negligible at atomic and subatomic scale. \n\nDespite the great success of the SM, one predicted particle has not yet been observed: the \\emph{Higgs boson}, an essential element for the inclusion of the mass of the particles in the equation of motion~\\cite{higgs1964,englert1964}. Because of this, a considerable effort is ongoing to prove or disprove the existence of the Higgs boson\\footnote{Some hints of its existence are confirmed by the present experiments~\\cite{tev_cmb, atlas_cmb, cms_cmb}.} and the analysis of the data collected by the CDF II experiment, situated at the Tevatron $p\\bar{p}$ collider, plays a relevant role in it~\\cite{cdf_cmb}.\n\nIn the scenario~\\cite{Tev4LHC} of a low-mass Higgs ($m_H\\lesssim 135$~GeV$\/c^{2}$), one of the most relevant CDF search channels is the $WH$ associate production with a $\\ell\\nu+b\\bar{b}$ final state. In light of this, the diboson decay channel considered for the presented analysis becomes a perfect benchmark for the Higgs boson search. The accurate SM prediction for the diboson production and decay can be used as a standard comparison for an unknown process. \n\nThis Chapter introduces the relevant aspects of the SM theory (Section\\ref{sec:sm_teo}), some of the latest results of the Higgs boson search (Section\\ref{sec:higgs_search}) and several, diboson related, experimental confirmation of the SM validity (Section\\ref{sec:dib_res}).\n\n\n\n\n\\section{The Standard Model of Elementary Particles}\\label{sec:sm_teo}\n\nThe SM is defined by the language of mathematics and theoretical physics so that it can be used to produce accurate predictions that have to be verified by the experiments. \nIn this language a particle is defined by a local quantum field. If no interaction is present, the free field is described by only two quantum numbers, the spin and the mass; if interactions are presents, the {\\em Gauge symmetries} can elegantly describe them: new quantum numbers classify the type and the strength of force while new particles, force-mediator vector bosons, are used to propagate the interaction.\n\nThe fundamental building blocks of matter, observed up to now, are the spin-$1\/2$ fields (\\emph{fermionic}), named \\emph{quarks} and \\emph{leptons}, and the spin-$1$ ({\\em vector}) fields, named \\emph{gauge bosons}. The leptons are divided into three \\emph{generations}, or \\emph{families}, and are\ngrouped in a left weak isospin doublet\\footnote{See Section~\\ref{sec:SM} for \nthe explanation of the weak isospin quantum number.} and a right weak isospin singlet. Also quarks are divided into three \\emph{flavor families} but weak \nisospin classification mixes quark doublets of different families. Quarks are also subject to the strong interaction, described by the \\emph{color} quantum number. Finally the {\\em charge} quantum number is used, for both quarks and leptons, to describe the electromagnetic interaction. \n\nThe force mediators are $W^\\pm$, $Z^0$, $\\gamma$, that carry electroweak \nforce, and $g$ (\\emph{gluons}), which mediate strong interaction. A short \nsummary of the SM fundamental particles is reported in Figure~\\ref{fig:particles}.\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.7\\textwidth]{Theory\/Figs\/SM_Particles}\n\\caption[SM Particles and Their Properties]{Quarks, leptons and gauge bosons in Standard Model and some of their \ncharacteristics~\\cite{pdg2010}, for each particle the corresponding antiparticle exists.}\\label{fig:particles}\n\\end{center}\n\\end{figure}\n\n\n\\subsection{Gauge Theory Example: QED}\n\nThe Quantum Electro-Dynamic (QED) is a perfect example to explain the importance of gauge invariance.\n\nThe equation that describes free fermionic fields is the {\\em Dirac} Lagrangian:\n\\begin{equation}\\label{eq:dirac}\n\\mathscr{L}(x)=\\bar{\\psi}(x)(i\\gamma^{\\mu}\\partial_{\\mu}- m)\\psi(x),\n\\end{equation}\nwhere $\\psi$ is the Dirac field of mass $m$ and $\\gamma^{\\mu}$ are the Dirac's \nmatrices. Equation~\\ref{eq:dirac} satisfies the \\emph{global} $U(1)$ symmetry\ntransformation: \n\\begin{equation}\n\\psi(x)\\to e^{iQ\\alpha}\\psi(x),\n\\end{equation}\nwith the electric charge $Q$ and the space independent parameter $\\alpha$ ($x$ is a space-time 4-vector).\nThe \\emph{Noether theorem}~\\cite{qftbook_weinberg} states that when a symmetry \nappears in a Lagrangian there is a corresponding conserved current. In the case of the Dirac field:\n\\begin{equation}\n\\partial_{\\mu}j^{\\mu} = 0\\quad\\mathrm{with}\\quad j^{\\mu} = -Q\\bar{\\psi}\\gamma^{\\mu}\\psi,\n\\end{equation}\ndescribes the conservation of charge, i.e the time component of the current 4-vector $j^{\\mu}$, integrated over the space, is a constant.\n\nAn elegant way to introduce interaction in the free Lagrangian is to shift from\nthe global, i.e. space independent, $U(1)$ transformation to a \\emph{local} \n$U(1)$ transformation, i.e. with a space dependent parameter $\\alpha(x)$:\n\\begin{equation}\n\\psi(x)\\to e^{iQ\\alpha(x)}\\psi(x).\n\\end{equation}\nTo maintain the gauge invariance condition in the Lagrangian~\\ref{eq:dirac}, a \ncovariant derivative $D_{\\mu}$ is introduced:\n\\begin{equation}\\label{eq:em_derivative}\n\\partial_{\\mu}\\to D_{\\mu}=\\partial_{\\mu}+iQA_{\\mu}, \\qquad\nD_{\\mu}\\psi(x)\\to e^{iQ\\alpha(x)}D_{\\mu}\\psi(x),\n\\end{equation}\nwhere the new vector field $A_{\\mu}$ is defined to transform in the following manner:\n\\begin{equation}\nA_{\\mu}\\to A_{\\mu}-\\frac{1}{Q}\\partial_{\\mu}\\alpha(x)\\mathrm{.}\n\\end{equation}\nEquations~\\ref{eq:em_derivative} and~\\ref{eq:dirac} can be composed to give the final QED Lagrangian:\n\\begin{equation}\\label{eq:L_qed}\n\\mathscr{L}_{QED}=\\bar{\\psi}(x)(i\\gamma^{\\mu}D_{\\mu}- m)\\psi(x) - \\frac{1}{4}F_{\\mu\\nu}F^{\\mu\\nu},\n\\end{equation}\nwhere $F_{\\mu\\nu}\\equiv\\partial_{\\mu}A_{\\nu}-\\partial_{\\nu}A_{\\mu}$ is the covariant kinetic term of $A_{\\mu}$.\nThe Dirac equation of motion for a field $\\psi$ undergoing electromagnetic interaction is obtained by applying the Euler-Lagrange equation~\\cite{qftbook_weinberg} to the QED Lagrangian:\n\\begin{equation}\n(i\\gamma^{\\mu}\\partial_{\\mu}- m)\\psi(x) = Q\\gamma^{\\mu}A_{\\mu}\\psi(x)\\mathrm{,}\n\\end{equation}\nthe force is mediated by the massless vector field $A_{\\mu}$. A mass term in \nthe form $\\frac{1}{2}m^{2}A_{\\mu}A^{\\mu}$ would break apart gauge invariance of\nEquation~\\ref{eq:L_qed}, indeed this is consistent with zero mass of the photon.\n\n\n\n\\subsection{Standard Model Theory}\\label{sec:SM}\n\nThe leptonic sector of the SM\\footnote{Only electroweak interaction on the leptons is considered here to simplify the discussion.} is based on the gauge group:\n\\begin{equation}\n SU(2)\\otimes U(1)\\textrm{,} \n\\end{equation}\nwhere $SU(2)$ is the non-Abelian group used in the spin algebra, and $U(1)$ is the Abelian group equivalent to the one used in QED. \nThe quantum number arising from $SU(2)$ is the \\emph{weak isospin}, $\\vec{T}$, and the one arising from $U(1)$ is {\\em hypercharge}, $Y$. They are related to the observed charge of real particles, $Q$, by the \nthe Gell-Mann-Nishijima equation:\n\\begin{equation}\nQ = T_{3}+\\frac{Y}{2}\\mathrm{.}\n\\end{equation}\nwhere $T_{3}$ is the third component of weak isospin.\n\nElectroweak interaction can be explained with a simplified model containing only two spin $1\/2$, elementary, massless, fermions, $f$ and $f'$, such that \n$Q_{f}=Q_{f'}+1$ ($Q$ is the electric charge). Weak interaction is built from V-A currents, i.e. left and right components are \ndefined and collected into a left doublet field and into two right singlet fields:\n\\begin{equation}\\label{eq:smL}\n\\psi_{1} \\equiv \\left( \\begin{array}{c} f_{L}(x) \\\\ f'_{L}(x) \\end{array}\\right), \\qquad \\psi_{2} \\equiv f_{R}(x) \n\\qquad \\psi_{3} \\equiv f'_{R}(x),\n\\end{equation}\nwith:\n\\begin{equation}\nf_{L,R}(x)=\\frac{1}{2}(1\\pm\\gamma_{5})f(x), \\qquad \n\\bar{f}_{L,R}(x)=\\frac{1}{2}\\bar{f}(x)(1\\pm\\gamma_{5}),\n\\end{equation}\n\\begin{equation}\nf'_{L,R}(x)=\\frac{1}{2}(1 \\pm \\gamma_{5})f'(x), \\qquad \n\\bar{f'}_{L,R}(x)=\\frac{1}{2}\\bar{f'}(x)(1 \\pm \\gamma_{5}).\n\\end{equation}\nAll the leptonic sector of the SM is explained by such pattern: the left doublet with $T_{3}=\\pm 1\/2$, $Y = 1$ is the charged lepton $f$\nplus the corresponding neutrino $f'$, while the right singlet with $T_{3}=0$, $Y = -2$ is only the charged lepton.\n\nThe electroweak interaction is introduced through $SU(2)\\otimes U(1)$ gauge transformation:\n\\begin{equation}\n\\psi_{j}(x)\\to\\psi'_{j}(x)= e^{i\\frac{\\tau}{2}\\cdot\\vec{\\alpha}(x) + iY_{j}\\beta(x)}\\psi_{j}(x),\n\\end{equation}\nof the free field Lagrangian:\n\\begin{equation}\\label{eq:freeL}\n\\mathscr{L_{0}}(x)=\\sum_{j=1}^{3}i\\bar{\\psi}_{j}(x)\\gamma^{\\mu}\\partial_{\\mu}\\psi_{j}(x)\\textrm{,}\n\\end{equation}\nwhere a covariant derivative is also introduced to maintain gauge invariance. The result is:\n\\begin{equation}\\label{eq:SML}\n\\mathscr{L}_{I}(x)=\\sum_{j=1}^{3}i\\bar{\\psi}_{j}(x)\\gamma^{\\mu}D_{\\mu}^{j}\\psi(x)_{j} - \\frac{1}{4}\\vec{W}_{\\mu,\\nu} \\vec{W}^{\\mu,\\nu} - \\frac{1}{4} B_{\\mu,\\nu}B^{\\mu,\\nu},\n\\end{equation}\n\\begin{equation}\\label{eq:bosons}\n\\mathrm{with}\\qquad D_{\\mu}^{j}= \\partial_{\\mu}-ig\\frac{\\tau}{2}\\cdot\\vec{W}_{\\mu}(x)-ig'Y_{j}B_{\\mu}(x)\\mathrm{,}\n\\end{equation}\n\\begin{equation}\\label{eq:tensor}\n\\mathrm{and}\\quad \\vec{W}_{\\mu,\\nu}=\\partial_{\\mu}\\vec{W}_{\\nu} - \\partial_{\\nu} \\vec{W}_{\\mu} + g \\vec{W}_{\\mu}\\times\\vec{W}_{\\nu}\\textrm{,}\\quad B_{\\mu,\\nu}= \\partial_{\\mu}B_{\\nu} - \\partial_{\\nu} B_{\\mu}.\n\\end{equation}\nEquation~\\ref{eq:bosons} contains three vector bosons ($\\vec{W}_{\\mu}$) from the $SU(2)$ \ngenerators, one vector boson ($B_{\\mu}$) from the $U(1)$ generator and four coupling constants: \n\\begin{equation}\n g,\\quad g'Y_{j} \\quad \\mathrm{with}\\quad j=1,2,3\\mathrm{.}\n\\end{equation}\nAfter some algebra the Lagrangian~\\ref{eq:SML} can be written in the form:\n\\begin{equation}\n\\mathscr{L}_{I}(x)= \\mathscr{L}_{CC}(x) + \\mathscr{L}_{NC}(x)\\mathrm{,}\n\\end{equation}\nwith a {\\em charged current} contribution ($\\mathscr{L}_{CC}$) and a {\\em neutral\ncurrent contribution} ($\\mathscr{L}_{NC}$). The charged current contribution is seen only by left doublet fields:\n\\begin{equation}\\label{eq:W}\n\\mathscr{L}_{CC}(x)=\\frac{g}{2\\sqrt{2}}\\Big\\{\\bar{f}(x)\\gamma^{\\mu}(1-\\gamma_{5})f'(x)\\frac{1}{\\sqrt{2}}W_{\\mu}^{+}(x) + h.c.\\Big\\},\n\\end{equation}\nwith $W_{\\mu}^{+}(x)$ defined by a linear combination of $W_{\\mu}^{1}(x)$ and $W_{\\mu}^{2}(x)$. Equation~\\ref{eq:W} defines the Lagrangian for charged current \ninteractions mediated by the $W$ boson.\n\nThe fermion coupling to $Z^{0}$ field and photon ($A$) field is produced in a similar way, by an appropriate orthogonal linear combination of neutral vector fields $B_{\\mu}(x)$\nand $W_{\\mu}^{0}(x)$:\n\\begin{equation}\n\\mathscr{L}_{NC}(x)= \\mathscr{L}_{NC}^{A}(x) + \\mathscr{L}_{NC}^{Z}(x)\\textrm{,}\n\\end{equation}\nwhere:\n\\begin{eqnarray}\n\\mathscr{L}_{NC}^{A}(x)=\\sum^{3}_{j=1}\\bar{\\psi}_{j}(x)\\gamma^{\\mu}\\big[g\\frac{\\tau_{3}}{2}\\sin \\theta_{W}\n+ g'Y_{j}\\cos \\theta_{W}\\big]\\psi_{j}(x)A_{\\mu}(x),\\\\\n\\mathscr{L}_{NC}^{Z}(x)=\\sum^{3}_{j=1}\\bar{\\psi}_{j}(x)\\gamma^{\\mu}\\big[g\\frac{\\tau_{3}}{2}\\cos \\theta_{W}\n+ g'Y_{j}\\sin \\theta_{W}\\big]\\psi_{j}(x)Z_{\\mu}(x),\n\\end{eqnarray}\nthe parameter $\\theta_{W}$ is named Weinberg angle and the generic four coupling constants, arising from the SM group structure, have now a physical meaning:\n\\begin{equation}\ng\\sin \\theta_{W} = e,\n\\end{equation}\n\\begin{equation}\ng'\\cos \\theta_{W}Y_{1}=e(Q_{f}-1\/2), \\qquad g'\\cos \\theta_{W}Y_{2}=eQ_{f},\n\\end{equation}\n\\begin{equation}\ng'\\cos \\theta_{W}Y_{3}=eQ_{f'}\\textrm{.}\n\\end{equation}\nPrevious equations are the core of the Standard Model. However one problem remains as no mass term appears for any of the fields: the spontaneous symmetry breaking and Higgs mechanism can generate the mass term without breaking the gauge invariance.\n\n\n\\subsection{Spontaneous Symmetry Breaking}\n\nSpontaneous symmetry breaking can be applied to Equation~\\ref{eq:smL} to give mass to $W^{\\pm}$ and $Z^{0}$ bosons. The actual application procedure is named {\\em Higgs mechanism}: two complex scalar fields are introduced such that they form an iso-doublet with respect to $SU(2)$:\n\\begin{equation}\\label{eq:phi_doublet}\n\\phi(x)\\equiv\\left(\\begin{array}{c} \\phi^{+}(x) \\\\\\phi^{0}(x) \\end{array}\\right),\n\\end{equation}\nthe field $\\phi^{+}(x)$ is the charged component of the doublet and $\\phi^{0}(x)$ is neutral component. The Higgs potential, $ V_{H}(x)$, is then defined as:\n\\begin{equation}\\label{eq:h_pot}\n V_{H}(x)\\equiv -\\mu^{2}\\phi^{\\dagger}(x)\\phi(x)-h\\big[\\phi^{\\dagger}(x)\\phi(x)\\big] ^{2},\n\\end{equation}\nwith $h>0$ and $\\mu^{2}<0$. The neutral scalar field $\\phi^{0}(x)$ has an unconstrained (i.e. to be obtained from measurements) vacuum expectation value of $\\frac{\\lambda}{\\sqrt{2}}$, so that (at first order)\nthe field $\\phi(x)$ is:\n\\begin{equation}\\label{eq:doublet}\n\\phi(x)=e^{\\frac{i}{\\lambda}\\vec{\\tau}\\cdot\\vec{\\theta(x)}}\\left(\\begin{array}{c} 0 \\\\\\frac{1}{\\sqrt{2}}\\big(\\lambda + \\chi(x)\\big) \\end{array}\\right),\n\\end{equation}\nwhere the $SU(2)$ gauge freedom is explicit. This permits to \\emph{gauge away} three of the four components of field $\\phi(x)$ leaving only one real scalar field: \n\\begin{equation}\\label{eq:}\n \\phi^{0}(x)= \\frac{1}{\\sqrt{2}}\\big(\\lambda + \\chi(x)\\big).\n\\end{equation}\nThe explicit evaluation of Equation~\\ref{eq:h_pot} and the coupling of $\\phi^{0}(x)$ with the electroweak force carriers ($W^{\\pm}$, $Z^0$) gives the last piece of the SM Lagrangian:\n\\begin{eqnarray}\\label{eq:finalL}\n\\mathscr{L}(x)&=& \\frac{1}{4}g^{2}\\lambda^{2}W_{\\mu}^{\\dagger}(x)W^{\\mu}(x) + \\frac{1}{1}(g^{2}+g'^{2})\\lambda^{2}Z_{\\mu}(x)Z^{\\mu}\\\\\n&&+\\frac{1}{2}g^{2}\\lambda W^{\\dagger}_{\\mu}(x)W^{\\mu}(x)\\chi(x)+\\frac{1}{4}g^{2}W^{\\dagger}_{\\mu}W^{\\mu}\\chi^{2}(x)\\nonumber\\\\\n&&+\\frac{1}{4}(g^{2}+g'^{2})\\lambda Z_{\\mu}(x)Z^{\\mu}(x)\\chi(x)+\\frac{1}{8}g^{2}Z_{\\mu}(x)Z^{\\mu}(x)\\chi^{2}(x)\\nonumber\\\\\n&&+\\frac{1}{2}\\big[\\partial^{\\mu}\\chi(x)\\partial_{\\mu}\\chi(x) + 2\\mu^{2}\\chi^{2}(x)\\big]\\nonumber\\\\\n&&+\\frac{\\mu^{2}}{\\lambda}\\chi^{3}(x)+\\frac{\\mu^{2}}{4\\lambda^{2}}\\chi^{4}(x)-\\frac{1}{4}\\lambda^{2}\\mu^{2}\\nonumber\\textrm{.}\n\\end{eqnarray}\nWe conclude that the $Z^{0}$ and $W^{\\pm}$ bosons have acquired mass:\n\\begin{equation}\n M_{W}=\\frac{1}{2}\\lambda g, \n\\end{equation}\n\\begin{equation}\n M_{Z}=\\frac{1}{2}\\lambda \\sqrt{g^{x}+g'^{x}}=\\frac{1}{2}\\frac{\\lambda g}{\\cos \\theta_{w} }, \n\\end{equation}\nsome parameters are now constrained, for example:\n\\begin{equation}\n M_{Z}=\\frac{M_{W}}{cos\\theta{w}}\\geqslant M_{W},\n\\end{equation}\n\\begin{equation}\n \\frac{G_{F}}{\\sqrt{2}}=\\frac{g^{2}}{8M^{2}_{W}},\n\\end{equation}\nwhile the Higgs mass, $M_{\\chi}=\\sqrt{-2\\mu^{2}}$ ($m_{H}$ is also used), remains a free parameter to be measured by the experiments. \nThe Higgs mechanism can generate also fermion masses if a Yukawa coupling is added:\n\\begin{eqnarray}\\label{eq:h_yukawa}\n\\mathscr{L_{f}}(x)&=&c_{f'}\\Bigg[(\\bar{f}(x),\\bar{f}'(x))_{L} \\left(\\begin{array}{c} \\phi^{+}(x) \\\\\\phi^{0}(x) \\end{array}\\right)\\Bigg]f_{R}'(x)\\\\\n&&+ c_{f}\\Bigg[(\\bar{f}(x),\\bar{f}'(x))_{L} \\left(\\begin{array}{c} -\\bar{\\phi}^{0}(x) \\\\\\phi^{-}(x) \\end{array}\\right)\\Bigg]f_{R}(x) + h.c.\\mathrm{,}\\nonumber\n\\end{eqnarray}\ntherefore, after symmetry breaking, fermion masses have the form:\n\\begin{equation}\nm_{f}=-c_{f}\\frac{\\lambda}{\\sqrt{2}}\\mathrm{,}\\qquad m_{f'}=-c_{f'}\\frac{\\lambda}{\\sqrt{2}}\\mathrm{,}\n\\end{equation}\nwhere the constants $c_{f}$ and $c_{f'}$ can be derived by the measurements of the fermion masses.\n\n\n\n\n\\section{Higgs Boson Search and Results}\\label{sec:higgs_search}\n\nThe mechanism that generates the mass of all the SM particles is a key element for the understanding of Nature, therefore it is not a surprise that the Higgs boson search is considered, by the High Energy Physics community, one of the most interesting research topics. \n\nAlthough the existence of the Higgs particle is unknown, its hypothetical couplings and decay properties are important for the interpretation of the experimental results: Figure~\\ref{fig:higgs_cx} shows the Higgs production cross section~\\cite{Tev4LHC}, at $\\sqrt{s}=1.96$~TeV and $\\sqrt{s}=7$~TeV, and Figure~\\ref{fig:higgs_br} shows the Higgs decay Branching Ratios~\\cite{LHCHiggs_cx_wg:2011ti} (BR) for a mass range $100< m_H<200$~GeV$\/c^{2}$. \n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[angle=-90,width=0.495\\textwidth]{Theory\/Figs\/tev_higgs_cx}\n\\includegraphics[angle=-90,width=0.495\\textwidth]{Theory\/Figs\/lhc_higgs_cx}\n\\caption[Higgs Boson Production Cross Sections]{Higgs boson production cross sections in different modes~\\cite{Tev4LHC}. At the Tevatron $p\\bar{p}$, $\\sqrt{s}=1.96$~TeV (left) and at the LHC $pp$, $\\sqrt{s}=7$~TeV (right). In 2012 the LHC raised the collision energy to $\\sqrt{s}=8$~TeV increasing still more the Higgs production cross section.}\\label{fig:higgs_cx}\n\\end{center}\n\\end{figure}\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.75\\textwidth]{Theory\/Figs\/YRHXS2_BR_Fig1-2}\n\\caption[Higgs Boson Decay Branching Ratios]{Higgs boson decay branching ratios~\\cite{LHCHiggs_cx_wg:2011ti}. The Higgs boson couples to the mass of the particles therefore the decay to the $b\\bar{b}$ quark pair is favored for $m_H\\lesssim 135$~GeV$\/c^{2}$ while $H\\to W^+W^-$ decay dominates for larger masses.}\\label{fig:higgs_br}\n\\end{center}\n\\end{figure}\n\nThe LEP experiments were the first to test the existence of Higgs boson for masses larger that $100$~GeV$\/c^{2}$, but, as no signal evidence was found~\\cite{lep}, all the searches were combined to provide a lower mass limit of $m_H>114.4$~GeV$\/c^{2}$, at 95\\% Confidence Level (CL). In the latest years also the experiments situated at the Tevatron and LHC colliders provided several mass exclusion limits~\\cite{tev_cmb, cms_cmb, atlas_cmb}. Figure~\\ref{fig:tev_cmb} gives a summary of the 95\\% CLs of all the three colliders, overlaid to the Tevatron result in the mass range $100< m_H<200$~GeV$\/c^{2}$, only a tiny fraction of the phase space is still available to the Higgs presence and, interestingly, a broad excess appears in mass range $110\\lesssim m_H \\lesssim 140$~GeV$\/c^{2}$. \n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.9\\textwidth]{Theory\/Figs\/tev_lhx_excl}\n\\caption[Tevatron Combined 95\\% Exclusion CL]{Mass exclusion limits (95\\% CL) obtained from the combination of all the Tevatron searches for a SM Higgs boson~\\cite{tev_cmb}. Exclusion limits obtained from the CMS~\\cite{cms_cmb} and Atlas~\\cite{atlas_cmb} Collaboration are overlaid.}\\label{fig:tev_cmb}\n\\end{center}\n\\end{figure}\n\nIn a short time, as the LHC continues the data taking, a conclusive statement about the Higgs existence will be possible.\n\nHowever, the Tevatron and LHC results are also complementary because they investigate different couplings of the Higgs boson.\nThe LHC experiments base most of the low-mass ($m_H\\lesssim 135$~GeV$\/c^{2}$) sensitivity on the $H\\to\\gamma\\gamma$ final state. This channel offers an excellent mass resolution and background rejection although at the price of a very low BR (see Figure~\\ref{fig:higgs_br}). This is optimal for the higher background rate and Higgs production cross sections availables at a the LHC $pp$ collisions of energy $\\sqrt{s}=7$~TeV and $\\sqrt{s}=8$~TeV.\n\nThe Tevatron experiments rely more on the $H\\to b\\bar{b}$ final state, where the Higgs is produced in association with a vector boson ($WH$ and $ZH$ production). The lower production cross section is compensated by the larger BR (see Figures~\\ref{fig:higgs_cx} and~\\ref{fig:higgs_br}) while the presence a leptonic decay of the $W$ or $Z$ boson allows to keep the background under control. Figure~\\ref{fig:tev_bb} shows the $H\\to b\\bar{b}$ only Tevatron combined search result. Furthermore the investigation of the $H\\to b\\bar{b}$ BR is important to understand the coupling of the Higgs with the fermion masses and for the confirmation of the SM assumption coming from Equation~\\ref{eq:h_yukawa}.\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.9\\textwidth]{Theory\/Figs\/tev_bb_cmb}\n\\caption[Tevatron $H\\to b\\bar{b}$ Combined 95\\% Exclusion CLs]{Mass exclusion limits (95\\% CL) obtained from the combination of the Tevatron searches~\\cite{tev_cmb} exploiting the $H\\to b\\bar{b}$ final state.}\\label{fig:tev_bb}\n\\end{center}\n\\end{figure}\n\n\n\n\n\n\\section{Status of the Diboson Measurements}\\label{sec:dib_res}\n\nIn the context of the Higgs searches at the Tevatron, the diboson observation in $\\ell\\nu + b\\bar{b}$ final state is particularly relevant as it is a direct check of the $p\\bar{p}\\to WH\\to \\ell\\nu + b\\bar{b}$ analyses.\n\nThe relevant tree-level diagrams involved in dibosons production are shown in Figure~\\ref{fig:wz_prod}. \n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.75\\textwidth]{Theory\/Figs\/wz_prod}\n\\caption[$WW$ and $WZ$ Tree-Level Diagrams]{$WW$ and $WZ$ production Feynman diagrams at tree-level, $t$-channel\n(left) and $s$-channel (right).}\\label{fig:wz_prod}\n\\end{center}\n\\end{figure}\nThe simultaneous emission of a $W$ and the $Z$ vector bosons can happen in the $t$-channel (left of Figure~\\ref{fig:wz_prod}), with the exchange of a virtual quark, or in the $s$-channel (right of Figure~\\ref{fig:wz_prod}) with the exchange of a virtual force carrier. The second case is due to the {\\em non-Abelian} characterisctics of the $SU(2)$ group that origins a {\\em Triple Gauge Coupling} (TGC) in the kinetic term (Equation~\\ref{eq:tensor}) of the SM Lagragnian. The cross sections for the $WW$ and $WZ$ production\\footnote{In this analysis we consider also $ZZ\\to \\ell\\ell + HF$ as a signal when a lepton is misidentified, however this contributes to less than 3\\% of the total diboson signal yield.} calculated at NLO~\\cite{mcfm_prd,mcfm_url}, for $p\\bar{p}$ collision at $\\sqrt{s}=1.96$~TeV, are:\n\n\\begin{equation}\n \\sigma_{p\\bar{p}\\to WW} = 11.34 \\pm 0.66\\mathrm{~pb};\\quad\\sigma_{p\\bar{p}\\to WZ} = 3.47 \\pm 0.21\\mathrm{~pb}; \n\\end{equation}\nAn increase in the TGC, $s$-channel, production cross section would point to a possible contribution from New Physics (NP) processes. However, the precision that we can obtain in the $\\ell\\nu+HF$ final state is not comparable to the one achievable in other channels with higher leptonic multiplicity (see Table~\\ref{tab:dib_measurements}).\n\nFigure~\\ref{fig:torte} shows the small dibosons BR in $b$ or $c$ quarks\\footnote{The experimental identification of $HF$ quarks, described in Section~\\ref{sec:secvtx} has also a low efficiency.}, furthermore \nthe hadronic final state is background rich and has a low invariant mass resolution. The search is challenging but it is a perfect standard candle to confirm the understanding of the $\\ell\\nu+HF$ dataset on a well known SM process.\n\n\\begin{figure}[!ht]\n\\begin{center}\n\\includegraphics[width=0.75\\textwidth]{Theory\/Figs\/wzpie}\\vspace{0.2cm}\n\\caption[Diboson Final States BR]{Branching ratios into the different final states allowed to $WW$ and $WZ$ production. $\\ell\\nu+ HF$ final state is highlighted in red.}\\label{fig:torte}\n\\end{center}\n\\end{figure}\n\n\nDiboson related experimental results are widely present in literature. LEP~\\cite{lep_diboson} performed the first measurements exploiting all the decay channels of the $WW$ and $ZZ$ processes: their cross sections were measured with good precision. The use of an $e^+e^-$ machine allowed also the observation of hadronically decaying $W$'s, in addition to the semi-leptonic $WW$ decays. However $WZ$ production was not allowed at LEP, since it is forbidden by charge conservation.\n\nHadron colliders, both Tevatron and LHC, observed $WW$, $WZ$ and $ZZ$ production in their fully leptonic decay modes, obtaining excellent agreement with the SM prediction~\\cite{cdfWZlllv_2012, d0WZZZlep_2012, cmsDib_2011, cdfZZllll_2012, atlasZZ_2012, cdfWWllvv_2012, d0WWllvv_2009, cmsWW_2012, atlasWW_2012, cdfDibMetJet_2009, cdfDiblvJet_Mjj2010, cdfDiblvJet_ME2010, d0DiblvJet_2011, cdfDiblvHF_Mjj2011, cdfWZZZCombHF_2012, d0WZZZCombHF_2011, tevWZZZCombHF_2012}.\n\nThe semi-leptonic final states, more difficoult to isolate due to the background rich hadronic environment, were observed at the Tevatron, both at CDF and D0 experiments. In this thesis we present the measurement in the channel $p\\bar{p}\\to WW\/WZ\\to \\ell\\nu +HF$ with an update and improvement of the analysis described in Appendix~\\ref{App:7.5}~\\cite{cdfDiblvHF_Mjj2011}, performed on a smaller dataset in 2011 ($7.5$~fb$^{-1}$). \n\nRecently~\\cite{cdfWZZZCombHF_2012, d0WZZZCombHF_2011, tevWZZZCombHF_2012}, CDF and D0 produced the evidence for $WZ\/ZZ$ production with $HF$ jets in the final state and $WW$ production considered as background. Both experiments produced single results in the three semileptonic diboson decay modes:\n\\begin{equation}\n WZ\/ZZ \\to \\ell\\ell+ HF\\textrm{;}\n\\end{equation}\n\\begin{equation}\n WZ\/ZZ \\to \\ell\\nu + HF\\textrm{;}\n\\end{equation}\n\\begin{equation}\n WZ\/ZZ \\to \\nu\\nu+ HF\\textrm{;}\n\\end{equation}\nwhere the $\\nu$ may indicate a lepton failing the identification. The analyses were performed as an exact replica of corresponding Higgs searches in those channels, with the final signal discriminants re-optimized for $WZ\/ZZ$ extraction. Heavy usage of multivariate tecniques, as for example in the $HF$ selection strategy~\\cite{HOBIT, BDT_D0} and in the final signal--background discrimination, is the key of the impressive sensitivity of these analyses. The final combined cross section measurement, with a significance of $4.6\\sigma$, is: \n\\begin{equation}\n \\sigma_{WZ+ZZ} = 4.47\\pm 0.67^{+0.73}_{-0.72}\\textrm{~pb,}\n\\end{equation}\nwhere the SM ratio between $WZ$ and $ZZ$ is imposed. This confirms the SM production prediction\\footnote{Both $\\gamma$ and $Z$ components are assumed in the neutral current exchange and corresponding production of dilepton final states for $75\\le m_{\\ell^+\\ell^{-}} \\le 105$~~GeV$\/c^{2}$.} of $\\sigma_{WZ+ZZ}=4.4\\pm0.3$~pb~\\cite{mcfm_prd,mcfm_url}. The evidence of the $WZ\/ZZ$ signal, obtained independently by each experiment and their combination, strongly supports the Tevatron $H\\to b\\bar{b}$ search results.\n\nA summary of all the present diboson measuremens is reported in Table~\\ref{tab:dib_measurements}.\n\n\\renewcommand{\\tabcolsep}{3pt}\n\\begin{table}[h]\n\\begin{center}\n \\begin{small}\n\n\\begin{tabular}{ccccc}\\toprule\nChannel & Experiment & $\\mathscr{L}$ (fb$^{-1}$) & Measured $\\sigma$(pb) & Theory $\\sigma$(pb) \\\\\\midrule\n\n$WZ\\to \\ell\\ell\\ell\\nu$ & CDF II\\cite{cdfWZlllv_2012} & $7.1$ & $3.9\\pm 0.8$ & $3.47 \\pm 0.21$ \\\\\n & D0\\cite{d0WZZZlep_2012} & $8.6$ & $4.5^{+0.6}_{-0.7} $ & $3.47 \\pm 0.21$ \\\\\n & CMS\\cite{cmsDib_2011} & $1.1$ & $17.0 \\pm 2.4\\pm 1.5$ & $ 17.3^{+1.3}_{-0.8}$ \\\\\n & Atlas\\cite{cdfWZlllv_2012} & $$ & $20.5 \\pm ^{+3.2+1.7}_{-2.8-1.5}$ & $ 17.3^{+1.3}_{-0.8}$ \\\\\n\\midrule\n\n$ ZZ \\to \\ell\\ell\n+\\ell\\ell\/\\nu\\nu$ \n & CDF II\\cite{cdfZZllll_2012} & $6.1$ & $ 1.64^{+0.44}_ {-0.38}$ & $1.4\\pm 0.1$\\\\\n& D0\\cite{d0WZZZlep_2012} & $8.6$ & $ 1.44^{+0.35}_{-0.34}$ & $1.4\\pm 0.1$\\\\\n\\midrule\n\n $ZZ \\to \\ell\\ell\\ell\\ell$ & Atlas\\cite{atlasZZ_2012} & $1.02$ & $ 8.5^{+2.7+0.5}_ {-2.3-0.4}$ & $6.5^{+0.3}_{-0.2}$\\\\\n & CMS\\cite{cmsDib_2011} & $1.1$ & $ 3.8^{+1.5}_{-1.2}\\pm 0.3 $ & $6.5^{+0.3}_{-0.2}$\\\\\n\n\\midrule\n\n$WW\\to \\ell\\ell\\nu\\nu$ & CDF II\\cite{cdfWWllvv_2012} & $3.6$ & $12.1\\pm 0.9 ^{+1.6}_{-1.4}$ & $11.3 \\pm 0.7$ \\\\\n & D0\\cite{d0WWllvv_2009} & $1.0$ & $11.5\\pm 2.2$ & $11.3 \\pm 0.7$ \\\\\n & CMS\\cite{cmsWW_2012} & $4.92$ & $52.4\\pm 2.0\\pm 4.7$ & $ 47.0\\pm 2.0 $ \\\\\n & Atlas\\cite{atlasWW_2012} & $1.02$ & $54.4\\pm 4.0 \\pm 4.4$ & $ 47.0\\pm 2.0$ \\\\\n\n\\midrule\n$ Diboson\\to \\nu\\nu + jets$ & CDF II\\cite{cdfDibMetJet_2009} & $3.5$ & $18.0\\pm 2.8 \\pm 2.6$ & $16.8 \\pm 0.5$ \\\\\n\n\\midrule\n$ WW\/WZ\\to \\ell\\nu + jets$ & CDF II\\cite{cdfDiblvJet_Mjj2010} & $4.3$ & $18.1\\pm 3.3\\pm 2.5$ & $16.8 \\pm 0.5$ \\\\\n & CDF II\\cite{cdfDiblvJet_ME2010} & $4.6$ & $16.5^{+3.3}_{-3.0}$ & \\\\\n\n & D0\\cite{d0DiblvJet_2011} & $4.3$ & $19.6^{+3.2}_{-3.0}$ & $16.8 \\pm 0.5$ \\\\\n \n\\midrule\n\n$ WW\/WZ\\to \\ell\\nu + HF$ & CDF II\\cite{cdfDiblvHF_Mjj2011} & $7.5$ & $18.1^{+3.3}_{-6.7}$ & $16.8 \\pm 0.5$ \\\\\n\n\\midrule\n\n & CDF II\\cite{cdfWZZZCombHF_2012} & $9.45$ & $4.1^{+1.8}_{-1.3}$ & $4.4 \\pm 0.3$ \\\\\n$ WZ\/ZZ \\to \\ell\\ell\/\\ell\\nu\/\\nu\\nu+ HF$ & D0\\cite{d0WZZZCombHF_2011} & $8.4$ & $5.0 \\pm 1.0 ^{+1.3}_{-1.2}$ & $4.4 \\pm 0.3$ \\\\\n & D0,CDF II\\cite{ tevWZZZCombHF_2012} & $7.5 - 9.5$ & $4.47\\pm 0.67^{+0.73}_{-0.72}$ & $4.4 \\pm 0.3$ \\\\\n\n\\bottomrule\n\\end{tabular}\n\\caption[Recent Diboson Measurements]{Summary of the recent measurements of the diboson production cross section in leptonic and semi-leptonic final states. The reference to the individual measurements are reported in the table as well as the used integrated luminosity, the statistical uncertainty appears before the systematic uncertainties (when both are available), the theoretical predictions are calculated with NLO precision~\\cite{mcfm_prd,mcfm_url}.}\\label{tab:dib_measurements}\n \n \\end{small}\n\\end{center}\n\\end{table}\n\\renewcommand{\\tabcolsep}{6pt}\n\n\\chapter*{\\centering \\begin{large}Abstract\\end{large}}\n\n\n\\clearpage\n\\include{Introduction\/Introduction}\n\n\n\\mainmatter\n\\clearpage\n\\include{Theory\/Theory}\n\n\\clearpage\n\\include{Detector\/Detector}\n\n\\clearpage\n\\include{PhysObj\/PhysObj}\n\n\\clearpage\n\\include{Sel\/Sel}\n\n\\clearpage\n\\include{Bkg\/Bkg}\n\n\\clearpage\n\\include{StatRes\/StatRes}\n\n\\clearpage\n\\include{Conclusions\/Conclusions}\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\nThe current generation of short pulse lasers reach intensities beyond $\\mathrm{I} > 10^{20}\\ \\mathrm{W\/cm^2}$, and hold great potential for many applications in high energy density science such as high energy x-ray backlighters and particle beams\\cite{Note1}\\cite{Note2}. While lower intensity lasers couple their energy into the transverse motion of electrons, relativistic electrons can couple energy into longitudinal motion via the laser's magnetic field. However, as intensity increases on relativistic laser platforms, the long ns pre-pulse pedestal that arrives prior to the main interaction intensifies as well. This pre-pulse, generated by amplified spontaneous emission processes in the laser, typically has an intensity contrast ratio $\\mathrm{(I_{Main Beam}\/I_{Pre-Pulse})}$ of $10^{6}$. In high intensity interactions this results in a pre-pulse of high enough intensity to ionize material, creating an underdense pre-plasma that extends for hundreds of\\SI{}{\\ \\micro\\meter}. \n\nWhile pre-plasma has been seen as detrimental to electron acceleration for applications such as fast ignition due to filamentation and self focusing instabilities\\cite{MacPhee}\\cite{Sunahara}, it actually can be beneficial in other ways, enabling a number of mechanisms that can generate energetic electrons. Accelerating electrons is the primary way to couple the energy in a laser to a plasma and provides the basis for a wide range of phenomena and applications. The generation of large quantities of x-ray and energetic secondary particles, such as ions \\cite{Pomerantz}, neutrons\\cite{Schollmeier} and positrons\\cite{Chen} highly depend on relativistic electrons. The electron acceleration in these experiments occur in a regime where the plasma response time is short compared to the laser pulse. In these experiments the density profile evolves slowly relative to the laser and develops into a quasi-steady state, opposite of the regime necessary for processes such as wakefield acceleration. Therefore, understanding the parameters that control the quantity, energy and trajectory of these electrons is critical to properly understanding the production of other energetic particles from secondary interactions.\n\nA common element for electron acceleration mechanisms is the important role played by quasi-static transverse and longitudinal pre-plasma electric fields in enhancing the energy transfer from the laser pulse to pre-plasma electrons \\cite{Booshan1}-\\cite{Robinson}. These fields are relatively weak compared to the field of the laser pulse and are unable to directly transfer considerable energy to the electrons. However, they do change the phase between the oscillating electric field of the laser and the electron velocity, which can result in a net energy gain with each laser period. This is the essence of the mechanism called direct laser acceleration (DLA) that leads to acceleration of the so-called super-ponderomotive electrons in an extended pre-plasma. These electrons have a corresponding relativistic $\\gamma$-factor greatly exceeding the conventional estimate of $\\gamma \\approx a_0^2$, where $a_0$ is the normalized laser amplitude. Indeed it has been shown that these super-ponderomotive electrons can enhance the energy of target normal sheath accelerated (TNSA) protons for radiography purposes but requires carefully controlled conditions for the rise time and trajectory of the electron beam \\cite{Kim}. \n\nPrior experiments using pulses with durations of 150 fs \\cite{Peebles},400 fs \\cite{Tanaka}, 500 fs \\cite{Toshi} and 700 fs \\cite{Ma} suggest that appreciable quantities of super-ponderomotive electrons should be accelerated for pulse lengths 400 fs or longer. The goal of the experiment presented here was to fill a gap between theory and experiment by accelerating large quantities of super-ponderomotive electrons using a very high intensity Texas Petawatt (TPW) laser while varying the interaction pulse length. These super-ponderomotive electrons were observed with energy up to $150$ MeV or more. However, contrary to prior experiments and simulations they were only detected at the longest pulse length tested (600 fs). In this paper we report the experimental measurement of super-ponderomotive electrons accelerated by DLA and the significance that a quasi static magnetic field has on their trajectory.\n\n\\section{Experimental Setup and Initial Simulations}\n\n\\begin{figure}[b]\n\t\\includegraphics[width=\\columnwidth]{Exp-Setup3}\n\t\\caption{Setup for the experiment on the TPW laser. A long pulse beam (green) generates an underdense plasma on the surface of the target while the high intensity short pulse (red) accelerates electrons measured by the primary EPPS diagnostics.}\n\\end{figure}\n\nSimilar to previous experiments \\cite{Peebles}-\\cite{Willingale} a low intensity, 2 ns beam with a large focal spot was used to generate a controllable, uniform, underdense plasma. This beam was injected 4 ns prior to the arrival of a short pulse beam, which had a variable pulse length (150, 450, 600 fs) and energy (30-105 J). Intensity was kept nominally to $2-3 \\times 10^{20}\\ \\mathrm{W\/cm^2}$ for all pulse lengths, and the beam was focused to a \\SI{10}{\\ \\micro\\meter} diameter spot, incident on the target at a $21.8^{\\circ}$ angle. The intrinsic pre-pulse from amplified spontaneous emission of the high intensity beam was nearly negligible after recent upgrades to the TPW, which improved the pre-pulse intensity contrast to over $3 \\times 10^{10}$ \\cite{GaulContrast}. Both beams had a wavelength of 1057 nm. The high-contrast laser was incident on a 1 $\\mathrm{mm^2}$ \\SI{100}{\\ \\micro\\meter} thick planar foil target composed primarily of aluminum, with a buried \\SI{20}{\\ \\micro\\meter} thick copper layer for diagnostic purposes. This foil was attached to a bulk plastic block used to inhibit electron refluxing, preventing double counting of electrons in our diagnostics. The experimental setup with the positions of diagnostics are shown in Figure 1.\n\n\\begin{figure}[b]\n\t\\includegraphics[width=\\columnwidth]{Sims1-a-revised3}\n\t\\caption{Simulation results of a 150, 450 and 600 fs beam (traveling left to right) incident on a target with underdense plasma. Electrons with energy $>$ 100 MeV traveling in the forward direction are color coded by angle relative to the target normal. The deep purple and green sections represent areas with significant electromotive force in the transverse direction on a 100 MeV electron.}\n\\end{figure}\n\nDiagnostics were fielded to characterize hot electron temperature and trajectory. While bremsstrahlung spectrometers (BMXS) and copper K$\\alpha$ imaging (SCI) were used, their signals are more sensitive to lower temperature electrons. High energy electrons were measured with EPPSes, a set of calibrated magnetic electron spectrometers \\cite{EPPS}. The spectrometers had 1 x 2 mm entrance pinholes, a strong magnet to deflect charged particles and differentiate them by energy and charge and two image plate detectors. The energy range was 5-100 MeV with energy resolution scaling with energy, from 5 keV at low energies to 1 MeV at higher energies. Two spectrometers were placed in the chamber, the first 43 cm away from the interaction facing the rear surface normal of the target and the second was placed 55.5 cm away from the interaction, 3 degrees off the front surface normal of the target. Spectra taken from the EPPS were characterized by the Half Maximum Integrated Energy (HMIE) value, which represents the energy value where 50 percent of the total energy in the spectrum is contained, e.g. a HMIE of 10 MeV means 50\\% of the energy in the spectrum is contained in electrons with less than 10 MeV energy. Hence a higher HMIE value will result from a spectrum with a greater proportion of high energy electrons.\n\nInitial hydrodynamic simulations with 2D FLASH \\cite{FLASH} were conducted to estimate the density profile of the underdense plasma using the parameters of the long pulse beam in the experiment. The resulting density profile was approximated as a sum of exponential decays with a density of around $10^{19}\\ \\mathrm{cm}^{-3}$ \\SI{150}{\\ \\micro\\meter} away from the solid surface. This density profile was input into the 2D particle-in-cell (PIC) code EPOCH \\cite{EPOCH} and truncated \\SI{160}{\\ \\micro\\meter} away from the target surface. The simulation box size was \\SI{250}{\\ \\micro\\meter} long, \\SI{60}{\\ \\micro\\meter} wide, with 10 electron and 5 proton macro particles per cell. The resolution was 30 cells\/\\SI{}{\\\\\\micro\\meter} in the x direction and 15 cells\/\\SI{}{\\\\\\micro\\meter} in the y direction. A laser was injected into the box with an intensity of $3 \\times 10^{20}\\ \\mathrm{W\/cm^2}$, a 70 fs rise time and a \\SI{10}{\\ \\micro\\meter} diameter spot.\n\nIn Figure 2 we present initial snapshots simulations for pulse lengths of 150, 450 and 600 fs taken just after the center of the pulse impacts the target in time. With the 150 fs pulse no significant super-ponderomotive electrons were accelerated. For 450 fs super-ponderomotive electrons are accelerated and travel primarily along the laser trajectory, while at 600 fs these electrons are dispersed by fields that develop near the critical density. Based on these simulations we anticipated measurable quantities of super-ponderomotive electrons ($\\mathcal{E} >$ 60 MeV, $a_0$ = 15.6) using pulse lengths greater than 400 fs.\n\n\\section{Results and Discussion}\n\nOver 35 experimental shots were performed over the course of the campaign with time split evenly between the 3 pulse length settings (150, 450 and 600 fs); half of the total shots did not use the long pulse beam to serve as the no pre-plasma cases. We verified that no super-ponderomotive electrons were measured with the shortest pulse length. When comparing the spectra for the 150 fs pulse length case there was no difference in the HMIE values when the long pulse beam was included. Upon increasing the pulse length to 450 fs the results were largely the same and no significant departure was seen between shots with and without pre-plasma. This was a surprising result since prior experiments \\cite{Tanaka}\\cite{Toshi} and our simulations with similar pulse lengths exhibited super-ponderomotive electrons and a change in spectrum shape.\n\n\\begin{figure}[b]\n\t\\includegraphics[width=\\columnwidth]{Exp-EPPS1-d2}\n\t\\caption{(a): Raw data from EPPS1 comparing cases with and without underdense plasma. Background subtracted lineouts are taken of the data and convolved with the spectrometer's calibrated dispersion to produce electron spectra. \\\\\\hspace{\\textwidth} (b): Electron spectra leaving the rear of the target measured by EPPS1 in the calibrated region of the diagnostic. Two 600 fs shots with pre-plasma shown in black and red produced substantial quantities of super-ponderomotive electrons $>$ 60 MeV resulting in a dramatic change to the electron spectrum compared to the no pre-plasma case (blue). \\\\\\hspace{\\textwidth}(c): HMIE values taken for all shots at 600 fs. Shots from the above spectra are circled with their respective color. At 600 fs, 3 out of 9 shots with underdense plasma measured a significant quantity of super-ponderomotive electrons}\n\\end{figure}\n\nExtending the pulse length further to 600 fs yielded significant changes to the electron spectrum shown in Figure 3; 3 out of the 9 shots with the long pulse beam had differently shaped spectra than their counterparts. These spectra contained electrons so energetic that they were detected at the end of the spectrometer (beyond its known calibration) and had energy exceeding 150 MeV. These results contradict the initial expectations from simulations shown in Figure 2. At 450 fs in simulations copious quantities of high energy electrons are accelerated nearly \\SI{50}{\\ \\micro\\meter} from the target surface and move forward into the target, while in the experiment no super-ponderomotive electrons were measured. At 600 fs super-ponderomotive electrons were measured on experiment while in simulations most electrons were severely deflected and few super-ponderomotive electrons propagated directly forward toward the diagnostic.\n\n\\begin{figure}[t]\n\t\\includegraphics[width=\\columnwidth]{Sims2-revised2}\n\t\\caption{Identical simulation setup to Figure 2 with the laser at a 20 degree incidence angle on target. The electrons are color coded by angle bin with respect to the EPPS diagnostic placed in the target normal direction. At 450 fs electrons travel in the direction of the laser similar to those in Figure 2 (b), while at 600 fs the electrons are deflected upwards in the target normal direction. This is due to the development of a large asymmetry in the electromotive force.}\n\\end{figure}\n\nWhat could cause this large discrepancy between the simulations and experimental data? The cause can be partly attributed to the EPPS diagnostic itself. The EPPS collects electrons only near the target rear normal and the solid angle ($\\sim 10^{-5}$ steradian) is an extremely small fraction of the solid angle where electrons are seen to travel in the simulations. Another difference between simulation and experiment lies in the inclusion of the laser incidence angle. In our experiment the laser was incident at roughly a $\\mathrm{20^o}$ angle and the EPPS diagnostic was placed facing the target normal, while our simulations used a normal incidence beam. When conducting another simulation with the same parameters but with the laser incidence angle and diagnostic angles included, the trajectories of super-ponderomotive electrons were altered significantly.\n\nFigure 4 shows the results from the angled simulation taken at the same times as those in Figure 2. Initially the $>$ 100 MeV electrons propagate primarily in the laser direction, not in the direction of the diagnostic. Measuring the angle of the electrons with respect to the viewing angle of the EPPS shows that the EPPS would miss the majority of super-ponderomotive electrons at 450 fs. However, later in time, electrons were scattered weakly toward the target normal rather than in all directions. At 600 fs the bulk of high energy electrons were deflected in the direction of the target normal and EPPS diagnostic, resulting in a greater chance of detection, which is directly supported by the experimental results.\n\n\\begin{figure}[t]\n\t\\includegraphics[width=\\columnwidth]{Exp-EPPS2-b}\n\t\\caption{HMIE values for electrons measured from the front surface of the target. Higher energy electrons were measured when an underdense plasma was present, regardless of pulse length suggesting that electrons in the pre-plasma were accelerated by the reflected beam.}\n\\end{figure}\n\\begin{figure}[]\n\t\\includegraphics[width=\\columnwidth]{Sims3-4-revised2}\n\t\\caption{(a): A current density plot for the 20 degree laser incidence angle case taken at the same time as Figure 4 (c). The effect of the asymmetric magnetic field on incident hot electrons is gradual causing a smaller deflection. Arrows indicate the location and direction of the majority of current flow.\\\\\\hspace{\\textwidth} (b): A plot of a simulation with a 10 degree laser incidence angle taken earlier in time. \\\\\\hspace{\\textwidth} (c): Plot of the 10 degree incidence simulation at the same time as the top 20 degree image. The current accelerated by the reflected beam expands and encounters the incident current more directly than in the 20 degree case. The super-ponderomotive electrons see a sudden, larger increase in the deflecting magnetic field as the fields join together.}\n\\end{figure}\n\nUpon closer examination of electron deflection at these later times it is found that the electrons encounter a counter-propagating current of colder electrons. The effect of this reverse current can be seen by mapping the electromagnetic force on a forward going electron (seen in purple and green in Figs 2 and 4). The forward going electrons generate a self confining magnetic field based on their current. Extending from the target surface, a counter-propagating current generates an opposite field that causes the electrons to split upwards or downwards. As pulse length increases this opposing field gains strength and causes more severe deflections. This effect has the largest impact on super-ponderomotive electrons, as they are accelerated via DLA far from the target surface and travel through this field. Lower energy, ponderomotive scaling electrons are less affected since they are accelerated near the critical density, closer to the target surface and do not interact with these fields over an extended distance.\n\nThe development of counter-propagating current can be traced to electrons in the pre-plasma near the target surface that are accelerated by reflected laser light away from the critical surface. This effect has been characterized in 3D PIC simulations conducted by F. Perez \\textit{et al.}\\cite{Perez}, which show the development of strong magnetic fields due to counter-propagating current after 473 fs. In the normal incidence case the two currents directly oppose each other and incoming electrons are deflected in all directions; in the angled case the opposing fields develop at an angle with respect to the target surface. This angle causes the fields from the counter-propagating current to build upon the fields from the incident current rather than oppose them, leading to a large asymmetry and significant deflection toward target normal.\n\nExperimental evidence of pre-plasma electrons accelerated by reflected light was found in the data of the spectrometer facing the front surface of the target (EPPS2). Measuring the HMIE values for electrons traveling in this direction shows distinctly higher energies for nearly all shots with the long pulse beam, regardless of pulse length (Fig 5). However, as our simulations show, the field from this reflected current requires time to develop and only has enough strength to deflect significantly electrons after 500 fs.\n\n\\section{Summary and Future Considerations}\n\nIn summary, we found that in experiments no super-ponderomotive electrons were detected in the forward direction for $\\tau_L$ = 150 fs, nor for $\\tau_L$ = 450 fs, in contradiction to earlier normal-incidence simulations, and only sporadically for $\\tau_L=$ 600 fs. Simulations including the experimental incidence angle showed that the super-ponderomotive electrons in fact were generated as early as $450$ fs, but were likely traveling in the direction of the laser direction and not towards the diagnostic placed in the target normal direction. As fields from an opposing current increased later in time these electrons were eventually deflected preferentially in the direction of the diagnostic. These results raise the question, how sensitive are super-ponderomotive electron trajectories to changes in laser incidence angle? To address this, another simulation with identical parameters was conducted, which used a 10 degree incidence angle rather than 20. As seen in figure 6 (c), electrons in the small angle case were deflected \\textit{more} significantly than in either of the previous two cases (0 and 20 degree), with a large portion deflected away from the target entirely late in time. Due to the smaller incidence angle, the region where laser light reflects off the critical density surface becomes narrower causing the opposing current to become more concentrated. More importantly, halving the incidence angle causes the relative angle between the reflected and incident light to be quartered. The opposing current therefore has a more head on trajectory with the incident electrons leading to greater deflection angle unlike the ``glancing blow'' seen in the large angle case.\n\nThese results point to a few considerations for future experiments and simulations. Understanding the effect of laser incidence angle is crucial for setting up an experiment. Many high intensity laser facilities do not allow for direct normal laser incidence due to the potential damage caused by reflected light, which means that incidence angle effects are a factor that must be considered for experiments that have been or will be performed at these facilities. Even for experiments that use perfectly normal incidence, the effect of counter-propagating current can result in deflecting electrons away from the laser direction as shown in Figure 2. Electron trajectories change significantly when incidence angle is varied and must be taken into consideration when placing diagnostics and when accelerating protons via TNSA using non-relativistic beams \\cite{Kim}. Due to the nature of high energy electrons, diagnostics capable of measuring super-ponderomotive electrons use very small solid angle apertures. Measuring these directional, high energy electrons poses a challenge to the diagnostics community to develop new tools capable of providing spatial information for higher energy electrons.\n\n\n\\section*{Acknowledgments}\n\nExperimental time and resources were supported by the United States Department of Energy under contract DE-NA0001858 (DOE). Simulations were performed using the EPOCH code (developed under UK EPSRC Grants No. EP\/G054940\/1, No. EP\/G055165\/1, and No. EP\/G056803\/1) using HPC resources provided by the TACC at the University of Texas. The work by A. Arefiev ,CHEDS researchers, and the TPW itself were supported by NNSA Cooperative Agreement DE-NA0002008.\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} diff --git a/data_all_eng_slimpj/shuffled/split2/finalzzqlhu b/data_all_eng_slimpj/shuffled/split2/finalzzqlhu new file mode 100644 index 0000000000000000000000000000000000000000..8ffb7645aa1fd8feb1c23918ab92bd1ca019f16d --- /dev/null +++ b/data_all_eng_slimpj/shuffled/split2/finalzzqlhu @@ -0,0 +1,5 @@ +{"text":"\\section{Introduction}\nQuantum Key Distribution (QKD) allows two remote parties, Alice and Bob, to generate a secret shared string of bits, that can be used for symmetric cryptography or other cryptographic protocols.\nUnlike classical key distribution, whose security is based on the computational difficulty of solving certain classes of problems, the security of quantum key distribution comes from the impossibility for an eavesdropper (Eve) to acquire information about an exchanged state without perturbing it.\nThis goal is obtained by using non-orthogonal states, as in the first QKD protocol, introduced by Bennett and Brassard in 1984 (BB84), that uses four different states of two mutually unbiased bases \\cite{Bennett1984}.\nIn 1992, Bennett showed that two non-orthogonal states are sufficient for QKD \\cite{Bennett1992}.\nThis protocol, however, has the drawback that the secure key rate is strongly affected by losses: indeed, Eve can extract information by increasing the losses and performing the so called unambiguous state discrimination (USD) attack~ \\cite{Tamaki2003,Tamaki2004}.\nThe addition of a third state is sufficient to make the B92 protocol unconditionally secure independently from the noise in the quantum channel \\cite{Fung2006,luca09pra}: however, rates comparable to the BB84 protocols can be obtained only when four states are detected at the receiver.\n\nThe optimal three-state QKD protocol, introduced in 2000 by Phoenix-Barnett-Chefles (PBC00)~\\cite{Phoenix2000}, uses states that form an equilateral triangle in the X-Z plane of the Bloch sphere.\nThe symmetry of this protocol can indeed be exploited to obtain rates comparable with the BB84 protocol but requiring only three detectors instead of four.\nThis protocol, however, still requires the public exchange of part of the sifted key in order to estimate the QBER.\nAn improvement of this protocol, introduced by Renes in 2004 (R04)~\\cite{Renes2004}, estimates the error rate from the number of inconclusive events, thus allowing to use all conclusive ones for key extraction.\nThe unconditional security of the PBC00 has been demonstrated, in the asymptotic case, in 2005 \\cite{Boileau2005}, for a bit error rate of up to $\\unit[9.81]{\\%}$.\nThe security of the PBC00 has been demonstrated also in the case of finite key \\cite{Mafu2014}, using the framework introduced by Scarani and Renner \\cite{Scarani2008} and the postselection technique \\cite{Christandl2009}.\nTheir work, however, does not consider the security of the R04 protocol, therefore excluding one of the most interesting features of this class of symmetrical codes.\nDespite all these theoretical results, no experimental implementation of equiangular three state QKD protocols has been reported so far.\nThe measurement apparatus proposed in the original work implements the trine measurement using an interferometric setup at the receiver~\\cite{Phoenix2000}.\nThis scheme requires careful alignment \\cite{Clarke2001,Clarke2001-1} and is not assured to have the long term stability required by Quantum Key Distribution (a stability of about half an hour is reported in \\cite{Clarke2001-1}).\nA new experimental scheme implementing the trine measurement using only passive optical elements has been proposed in \\cite{Saunders2012}.\nUsing this apparatus, we demonstrate the feasibility of equiangular three state QKD and assess its key generation rate, both in the asymptotic limit of infinite key and taking into account finite key effects.\n\n\\section{Results}\n\\subsection{Protocol}\nThe R04 protocol uses three quantum states, $\\{ \\ket{\\psi_1}, \\ket{\\psi_2}, \\ket{\\psi_3} \\}$, placed in an equilateral triangle in the X-Z plane of the Bloch sphere, as shown in Figure \\ref{fig:bloch}.\n\\begin{figure}[t]\n \\centering\n \\includegraphics[width=0.7\\linewidth]{images\/bloch.pdf}\n \\includegraphics[width=0.7\\linewidth]{images\/povm.pdf}\n \\caption{States used in the R04 protocol (above) and POVM $\\{ \\Pi_i \\}$ used for the measurement (below). The states lie in the X-Z plane of the Bloch sphere. They are grouped into the sets $S_1 = \\{ \\ket{\\psi_1}, \\ket{\\psi_2} \\}$, $S_2 = \\{ \\ket{\\psi_2}, \\ket{\\psi_3} \\}$, and $S_3 = \\{ \\ket{\\psi_3}, \\ket{\\psi_1} \\}$, where the first element of each set corresponds to bit 0 and the second to bit 1. The POVM is implemented by using a partially polarizing beam-splitter (pPBS), a half-wave plate (HWP) at $22.5^\\circ$ and a polarizing beam-splitter (PBS).}\n \\label{fig:bloch}\n\\end{figure}\nThe states are grouped into three different sets, $S_1 = \\{ \\ket{\\psi_1}, \\ket{\\psi_2} \\}$, $S_2 = \\{ \\ket{\\psi_2}, \\ket{\\psi_3} \\}$, and $S_3 = \\{ \\ket{\\psi_3}, \\ket{\\psi_1} \\}$. \nIn each set, the first state is associated with the bit $0$ and the second with the bit $1$.\nDifferently from other QKD protocols, each state brings no information about its associated bit before the information about the used set is disclosed. \\\\\n\nWe implemented an entanglement-based version of the protocol, using polarization-entangled photon pairs in the singlet state\n\\begin{equation}\n\t\t\t\t\\ket{\\Psi^{-}} = \\frac{\\ket{H}_A \\ket{V}_B - \\ket{V}_A \\ket{H}_B}{\\sqrt{2}},\n\t\t\t\t\\label{eq:psi-}\n\\end{equation}\nwhere the subscripts $A$ and $B$ indicate the photons going to, respectively, Alice's and Bob's detection apparatus.\nIn this state, photons A and B are anti-correlated in any measurement basis.\nAlice measures photon $A$ using the POVM $\\{\\Pi_i\\equiv\\frac23\\ket{\\psi_i^\\perp}\\bra{\\psi_i^\\perp}\\}$, with $\\ket{\\psi_1^\\perp} = \\ket{V}$, $\\ket{\\psi_2^\\perp} = \\frac{\\sqrt{3}}{2}\\ket{H} - \\frac{1}{2}\\ket{V}$, $\\ket{\\psi_3^\\perp} = \\frac{\\sqrt{3}}{2}\\ket{H} + \\frac{1}{2}\\ket{V}$. In our notation the state $\\ket{\\psi_i^\\perp}$ is orthogonal to $\\ket{\\psi_i}$.\nWhen Alice obtains a detection in the state $\\ket{\\psi_i^\\perp}$ (each with probability $\\frac{1}{3}$), she is sending to Bob the state $\\ket{\\psi_i}$ where $\\ket{\\psi_1} = \\ket{H}$, $\\ket{\\psi_2} = \\frac{1}{2}\\ket{H} + \\frac{\\sqrt{3}}{2}\\ket{V}$, and $\\ket{\\psi_3} = \\frac{1}{2}\\ket{H} - \\frac{\\sqrt{3}}{2}\\ket{V}$.\nThis operation corresponds to the random preparation, with equal probability, of one of the three states $\\{\\ket{\\psi_1}, \\ket{\\psi_2}, \\ket{\\psi_3}\\}$, as in the prepare-and-measurement scheme of the R04 described in \\cite{Boileau2005,Scarani2009}.\nBob performs his measurements in the same POVM as Alice $\\{ \\Pi_i \\}$.\nAfter all measurements, Bob and Alice compare the\ninstants of their events, \nkeeping only those where both have a detection within a fixed coincidence window. \n\nEven if they have already exchanged all symbols, Alice and Bob do not share any bit string yet, because each state can mean both 0 or 1.\nAlice uses a QRNG to choose the bit value for each symbol.\nThe combination of the state and the bit value unambiguously determines the set $S_i$ used for that event (for example, if Alice sends $\\ket{\\psi_2}$ and the QRNG gives $1$, the set used for that event is $S_1$).\nFor each event, Alice tells Bob the corresponding set by sending him the value of the index $i$.\nBob uses $i$ to associate $\\ket{\\psi_2^\\perp}$ (for $i = 1$), $\\ket{\\psi_3^\\perp}$ (for $i = 2$), and $\\ket{\\psi_1^\\perp}$ (for $i = 3$) with bit 0, and $\\ket{\\psi_1^\\perp}$ (for $i = 1$), $\\ket{\\psi_2^\\perp}$ (for $i = 2$), and $\\ket{\\psi_3^\\perp}$ (for $i = 3$) with bit 1.\nAll other combinations are marked as inconclusive, since Bob is not able to determine the state sent by Alice.\nBob tells Alice which events are inconclusive and they both discard them.\nThey then estimate the quantum bit error rate (QBER) from the fraction of inconclusive events \\cite{Renes2004,Boileau2005}, and use this information to distill the key using error correction and privacy amplification \\cite{Scarani2008-1}. \n\n\\subsection{Setup}\nThe experimental setup is shown in Figure \\ref{fig:setup}.\n\\begin{figure}[t]\n \\centering\n \\includegraphics[width=\\linewidth]{images\/Sagnac_noPOVM.pdf}\n \\caption{Experimental setup used for the experiment. The light issued by a laser diode (LD) at $\\unit[404.5]{nm}$ is collected into a polarization maintaining single mode fiber (PM-SMF) for spatial mode filtering. The polarization is adjusted by a PBS followed by a HWP and the beam is focused by a $\\unit[200]{mm}$ lens (L) into the center of the interferometer. The interferometer is formed by a PBS, a dual-wavelength HWP (d-HWP) and a PPKTP crystal, where the down-conversion takes place. Photons exiting the interferometer are collected into single mode fibers (SMF), preceded by a long-pass filter (LPF) to filter out residual pump intensity. Fiber birefringence is compensated by two quarter-wave plates (QWP) and a half-wave plate (HWP). The receiver consists in an implementation of the POVM $\\{ \\Pi_i \\}$, with output Ai and Bi corresponding to $\\ket{\\psi_i^\\perp}\\bra{\\psi_i^\\perp}$ for, respectively, Alice and Bob. All the output signals are timed by a common timetagger.} \n \n \n\\label{fig:setup}\n\\end{figure}\nEntangled photon pairs are produced by using a $\\unit[30]{mm}$ periodically poled KTP crystal in a polarization-based Sagnac interferometer \\cite{Kim2006,Fedrizzi2007}.\nThe source is pumped with a continuous wave (CW) laser at $\\unit[404.5]{nm}$, with a power of $\\unit[3.5]{mW}$.\nThe down-converted photons have a central wavelength of $\\unit[809]{nm}$, with $\\unit[0.2]{nm}$ full width at half maximum (FWHM), and are collected into single-mode fibers.\nIn this configuration, the setup has a mean coincidence rate of $\\unit[29]{kHz}$, with a 5\\% heralding ratio.\nThe fraction of multi-pair over one pair events, measured by putting one output of the source into a Hanbury-Brown-Twiss interferometer, is $\\sim 3 \\cdot 10^{-3}$.\nAmong these events, only those which are partially correlated in the polarization degree of freedom are exploitable by Eve through the photon number splitting (PNS) attack \\cite{Scarani2008}.\nThe ratio of the number of these events to all multi-pair ones is $\\zeta \\simeq \\frac{\\tau_c}{\\Delta t} = 5 \\cdot 10^{-3}$, where $\\tau_c = \\unit[8]{ps}$ is the coherence time of down-converted photons and $\\Delta t = \\unit[1.5]{ns}$ is the coincidence window.\nThe fraction of correlated multi-photon events over the total number of detection events is $\\sim 1.5 \\cdot 10^{-5}$, thus the information leaked to Eve is negligible.\nA set of two quarter-wave plates (QWP) and one half-wave plate (HWP) is placed at the exit of the fiber at Bob's side in order to compensate polarization rotations induced by fiber birefringence. \n\nThe receiving apparatus implementing the POVM $\\{\\Pi_i\\}$ consists of a partially polarizing beam-splitter (pPBS), that completely transmits the horizontal polarization and has a reflectivity of $\\unit[66.7]{\\%}$ for the vertical polarization, followed by a HWP at $\\theta = 22.5^\\circ$ and a polarizing beam-splitter (PBS).\nGiven an arbitrary input state $\\ket{\\phi} = \\alpha \\ket{H} + \\beta \\ket{V}$, with $|\\alpha|^2 + |\\beta|^2 = 1$, the pPBS routes it to detector $1$ with probability $P_1 = \\frac{2}{3}|\\beta|^2$.\nThe state at the transmitting output of the pPBS, $\\alpha \\ket{H} + \\frac{1}{\\sqrt{3}} \\beta \\ket{V}$, is transformed into $\\alpha \\ket{-} + \\frac{1}{\\sqrt{3}} \\beta \\ket{+} = \\frac{1}{\\sqrt{2}} ( \\alpha + \\frac{1}{\\sqrt{3}} \\beta) \\ket{H} - \\frac{1}{\\sqrt{2}} ( \\alpha - \\frac{1}{\\sqrt{3}} \\beta) \\ket{V}$ by the\nHWP.\nThus, the probabilities of a click at detectors $2$ and $3$ are, respectively, $P_2 = \\frac{1}{2} |\\alpha - \\frac{1}{\\sqrt{3}} \\beta |^2$ and $P_3 = \\frac{1}{2} | \\alpha + \\frac{1}{\\sqrt{3}} |^2$.\nIt is easy to check that the above\nprobabilities can be also written as $P_i = \\bra{\\phi} \\Pi_i \\ket{\\phi}$, with $\\Pi_i = \\frac{2}{3} \\ket{\\psi_i^\\perp} \\bra{\\psi_i^\\perp}$: then, the above described apparatus implements the POVM $\\{\\Pi_i\\}$.\nPhotons are detected using silicon single photon counting modules (SPCM), characterized by a dead time of $\\unit[21]{ns}$ and a jitter of $\\sim800$ ps FWHM. Detection events are time-tagged with a resolution of $\\unit[81]{ps}$. \n\n\\subsection{Data acquisition}\nA two hour continuous run of the apparatus has led to the exchange of about $10^{9}$ symbols within a coincidence window of $\\unit[1.5]{ns}$.\nThe events can be described as pairs (Ai,Bj), with $i$ the number of the detector clicking at Alice's side and $j$ the one at Bob's side.\nEvent distribution is shown in table \\ref{tab:data} and Figure \\ref{fig:histogram}.\n\\begin{table}[h]\n \\centering\n \\begin{tabular}{|c|c|c|c|}\n \\hline\n & A1 & A2 & A3 \\\\\n \\hline\n B1 & 0.6 & 35.8 & 33.6 \\\\\n B2 & 35.1 & 0.6 & 32.8 \\\\\n B3 & 33.4 & 33.2 & 0.4 \\\\\n \\hline\n \\end{tabular} \n \\caption{\\label{tab:data} Total number of coincidences at the different detectors (million events). The cell (Ai,Bj) corresponds to a coincidence of Alice's detector $i$ and Bob's detector $j$.}\n\\end{table}\n\\begin{figure}[t]\n \\centering\n \\includegraphics[width=\\linewidth]{images\/histogram.pdf}\n \\caption{Total number of coincidences at the different detectors. Full (red) bars correspond to detected events and (blue) contours represent the expected number of detection events. }\n \\label{fig:histogram}\n\\end{figure}\nAfter the collection of all data, a QRNG \\cite{Vallone2014} is used to generate the bit value for each symbol.\nCoincidence events are then analyzed using the sifting procedure summarized in Table \\ref{tab:sifting}.\n\\begin{table}[h]\n \\centering\n \\begin{tabular}{|c|c|c|c||c|c|c|}\n \\hline\n & \\multicolumn{3}{c||}{Bit = 0} & \\multicolumn{3}{c|}{Bit = 1} \\\\\n \\hline\n & A1 & A2 & A3 & A1 & A2 & A3\\\\\n \\hline\n B1 & 1 & Inc & 0 & 0 & 1 & Inc \\\\\n B2 & 0 & 1 & Inc & Inc & 0 & 1 \\\\\n B3 & Inc & 0 & 1 & 1 & Inc & 0\\\\\n \\hline\n \\end{tabular}\n \\caption{\\label{tab:sifting} Sifting procedure, according to the random choice of the bit at Alice's side (on the left for 0, on the right for 1). The cell (Ai,Bj) stands for a coincidence between Alice's detector $i$ and Bob's detector $j$. Inconclusive events are marked as ``Inc''. The events in the diagonal (Ai,Bi) give an error independently from the bit choice. The other combinations (Ai,Bj), with $i \\neq j$, are either a ``good'' conclusive or an inconclusive event, according to Alice's choice.}\n\\end{table}\nThe events of the form (Ai,Bi) are bit errors, while the others are either a ``good'' conclusive or an inconclusive result according to Alice's choice.\nThe string of conclusive results gives the sifted key, from which a secret key can be distilled using classical post-processing. \n\n\\subsection{Secret key rate}\nPost-processing consists of a series of passages that transform a partially correlated, partially secret key into a new one Eve has negligible information of \\cite{Scarani2008-1,Scarani2009}.\nThe effect of these tasks is a reduction of the number of bits and can be quantified using the secret fraction $r$, defined as the ratio between secure and conclusive bits \\cite{Scarani2009}.\nIn the asymptotic limit of infinitely long key, the key fraction of the R04 is given by \\cite{Boileau2005}\n\\begin{equation}\n r = 1 - f_{EC} h(Q) - h\\left({\\frac{5}{4}Q}\\right),\n \\label{eq:asymp_r}\n\\end{equation}\nwhere $f_{EC} = 1.1$ is the efficiency of the error correction protocol \\cite{Mateo2015}.\nThe number of secure bits is given by $N_{conc} r$ and, dividing it by the exposure time, the secret key rate is obtained. \n\nFigure \\ref{fig:keyrate} shows the behavior of the QBER and of the key rate during a two hour acquisition.\nThe slight reduction in the sifted key rate is probably due to a misalignment in the fiber coupling of the entangled source.\nThe losses can be estimated from the ratio of coincidences over single counts.\nThe measured $5\\%$ heralding efficiency corresponds to a total loss level of $\\unit[13]{dB}$, with a contribution of $\\unit[1.5]{dB}$ due to the POVM.\n\\begin{figure}[t]\n \\centering\n \\includegraphics[width=\\linewidth]{images\/keyrate.pdf}\n\t \\caption{Result obtained during 2 hour of continue acquisition. The time has been divided into $90$ blocks of about $\\unit[80]{s}$ each, with a mean number of $1.1 \\cdot 10^{6}$ sifted bits. The QBER is estimated from inconclusive events, with Poissonian error bars. The secret key rate is estimated for the case of infinitely long key, using equation (\\ref{eq:asymp_r}).}\n \\label{fig:keyrate}\n\\end{figure}\nThe QBER is estimated as $Q = \\frac{1-2I}{1-I}$, where $I$ is the fraction of inconclusive results \\cite{Boileau2005}.\nThe QBER remains almost constant at a level below $2\\%$ for all the acquisition, thus confirming the stability of both the source and the POVM $\\{ \\Pi_i \\}$ during the acquisition.\nThe visibility of the source in two mutually unbiased bases at the exit of the Sagnac interferometer, before fiber injection, has been measured to be between $97\\%$ and $98\\%$: then the measured QBER level can be attributed almost completely to the source.\nA small contribution to the QBER is due to the small imbalances between the different channels of the POVM $\\{ \\Pi_i \\}$.\nThese values, estimated by computing the ratio between the raw counts at the different detectors, vary between $0.95$ and $1.05$, in line with what observed in the previous implementation of the POVM \\cite{Saunders2012}.\n\nIn a real scenario, the number of exchanged signals is always finite and the security analysis must take this fact into account.\nThe finite key analysis of the R04 protocol is very similar to the one of the PBC00 \\cite{Mafu2014}, the only substantial difference lying in the estimation of the bit error rate.\n\\begin{figure}[t]\n\\includegraphics[width=\\linewidth]{images\/finite.pdf}\n\\caption{Finite key analysis of the R04 protocol.\nThe y-axis represents the fraction of exchanged symbols giving a secure bit. Each point is calculated using the first N collected data for the estimation of the QBER. The security parameter for both collective and general attacks is fixed to $4 \\cdot 10^{-10}$.}\n\\label{fig:finite}\n\\end{figure}\nThe PBC00 estimates it by comparing part of the sifted key through the public channel, while the R04 use the fraction of inconclusive results.\nThe method is based on the fact that the choice between ``good'' conclusive and inconclusive results is given by a random event at Alice's side after the exchange of all the qubits, therefore Eve has no way of differentiate between the two and their number is approximately equal \\cite{Boileau2005}.\nDefining $I$ the fraction of inconclusive results, the fraction of ``good'' conclusive ones can be written as $(1-Q)(1-I)$.\nUsing the Hoeffding bound \\cite{Hoeffding1963}, the inequality\n\\begin{equation}\n\t\t\t\t\\left| (1-Q)(1-I) - I \\right| \\leq \\xi(\\epsilon,N) := \\sqrt{\\frac{2}{N}\\log\\frac{2}{\\epsilon}}\n\t\t\t\t\\label{eq:hoeffding}\n\\end{equation}\nis valid with probability at least $1-\\epsilon$.\nThis implies that, with the same probability, the QBER $Q$ is less than $\\widetilde{Q} := \\frac{1-2I+\\xi(\\epsilon,N)}{1-I}$.\nThe secret key fraction, in the case of collective attacks, is \\cite{Scarani2008,Mafu2014}\n\\begin{align}\n\t\t\t\tr_{col} = & \\left[ 1 - h\\left( \\frac{5}{4} \\widetilde{Q} \\right) \\right] - 7 \\sqrt{\\frac{1}{N} \\log_2{\\frac{2}{\\bar{\\epsilon}}}} \\nonumber \\\\\n\t\t\t\t & - \\frac{1}{N}\\log_2{\\frac{1}{\\epsilon_{EC}}} - \\frac{1}{N}\\log_2{\\frac{2}{\\epsilon_{PA}}} - f_{EC} h(Q)\n\t\t\t\t\\label{eq:r_finite}\n\\end{align}\nfor $N$ exchanged signals, with security parameter $\\epsilon_{col} := \\bar{\\epsilon} + \\epsilon_{EC} + \\epsilon_{PA} + \\epsilon_{PE}$.\nThis result can be extended to the case of general attacks by exploiting the postselection technique \\cite{Christandl2009}, giving \\cite{Mafu2014} \n\\begin{equation}\nr_{gen} = r_{col} - \\frac{6 \\log_2{(N+1)}}{N},\n\\label{eq:r_finite_gen}\n\\end{equation}\nwith security parameter $\\epsilon_{gen} = (N+1)^{3} \\epsilon_{col}$. \n\nFigure \\ref{fig:finite} shows the secret key fraction of the R04 protocol in the finite key scenario.\nFor each point, both the QBER and the number of conclusive events is evaluated on the first $N$ exchanged symbols.\nThe secret key fraction is calculated by using equations (\\ref{eq:r_finite}) and (\\ref{eq:r_finite_gen}).\nFor collective attacks, the chosen security parameter is $\\epsilon_{col} = 4 \\cdot 10^{-10}$, with $\\bar{\\epsilon} = \\epsilon_{EC} = \\epsilon_{PA} = \\epsilon_{PE} = 10^{-10}$.\nThe same value has been chosen for $\\epsilon_{gen}$, therefore the term $r_{col}$ of equation (\\ref{eq:r_finite_gen}) is calculated using $\\bar{\\epsilon} = \\epsilon_{EC} = \\epsilon_{PA} = \\epsilon_{PE} = \\frac{10^{-10}}{(N+1)^3}$.\nThe plots show that at least $10^{4}$ - $10^{5}$ signals are necessary to exchange a key, while already $N = 10^{6}$ (slightly more than half a minute at $\\unit[29]{kHz}$) gives a reasonable key fraction.\nThe difference between the key fraction for collective and general attacks is more marked for lower values of $N$ and tends to disappear for large number of exchanged symbols, where both approach the asymptotic key fraction. \n\n\\section{Discussion}\nIn this work, we demonstrated the experimental feasibility of the equiangular three state QKD protocol R04.\nWe showed that the scheme proposed by \\cite{Saunders2012} for the POVM is suitable for applications in Quantum Key Distribution.\nState preparation was simplified by using an entanglement-based version of the protocol, with the same POVM at both Alice's and Bob's side.\nWe also showed that the estimation of the bit error rate from inconclusive results is feasible in the finite key scenario.\nThe implemented scheme was demonstrated to be stable and highly reliable, allowing a two hour data acquisition without any significative change in the QBER value.\nThe performance of the protocol is comparable with the BB84, despite the less efficient parameter estimation, both in the asymptotic limit and for finite key.\nIts simpler receiving apparatus, requiring only three single photon detectors, makes it a valid alternative to current implementations of QKD based on the BB84 protocol.\nFinally, this work extends the experimental investigation of equiangular spherical codes to a still uncovered area of quantum information: quantum key distribution.\n\n\\acknowledgements\nWe thank Davide G. Marangon and Nicola Laurenti for helpful discussions.\nThis work has been carried out within the Strategic-Research-Project QUANTUMFUTURE of the University of Padova and the Strategic-Research-Project QUINTET of the Department of Information Engineering, University of Padova.\nM.S. acknowledges the Ministry of Research and the Center of Studies and Activities for Space (CISAS) ``Giuseppe Colombo'' for financial support.\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\\label{sec:intro}\n\nEleven-dimensional, \\ensuremath{\\mathfrak{osp} \\left( 32 \\middle\\vert 1 \\right)}-based Chern--Simons (CS) Supergravity was put forward by\nR.~Troncoso and J.~Zanelli (TZ) in 1997~\\cite{Tro97,Tro98}.\nThe appearance of this eleven-dimen\\-sional theory followed previous efforts in lower dimensions by\nS.~Deser, R.~Jackiw and S.~Templeton~\\cite{Des82b,Des83},\nA.~Ach\u00facarro and P.~K.~Townsend~\\cite{Ach86}, and\nA.~H.~Chamseddine~\\cite{Cha90}.\n\nThe ``supergravity'' name has proven a bit controversial.\nOn one hand, the TZ theory includes gravity and fermionic matter and sports an exact, off-shell\n\\ensuremath{\\mathfrak{osp} \\left( 32 \\middle\\vert 1 \\right)}\\ symmetry which mixes bosons and fermions.\nOn the other, it differs quite radically%\n\\footnote{For instance, the TZ theory lacks equality between bosonic and fermionic degrees of freedom (at least in the off-shell counting).\nThis matching, which has been taken as a watermark of supersymmetry in the literature, lies also at the foundations of the MSSM.\nThe recent negative results of the search for supersymmetric particles at the LHC~\\cite{CMS11a,Str11,ATL11a,ATL11b,Buc11},\nalthough clearly still at an early stage,\nshould not then be taken as an irrefutable sign that supersymmetry as a principle is wrong.}\nfrom the ``standard'' 1978 theory of Cremmer, Julia and Scherk (CJS)~\\cite{Cre78},\nwhich is widely regarded as unique~\\cite{Des97}.\n\nThe purpose of this paper is to shed new light on the relation between the TZ and the CJS theories.\n\nIn this regard, we would like to highlight here the work of\nM.~Ba\u00f1ados~\\cite{Banh01}, who showed that the linearized perturbations of the TZ theory are the same as the ones of CJS supergravity\n(see also Ref.~\\cite{Hor97}).\nThat result, however, is based on the identification of the CJS three-form $A_{3}$ with $e^{a} T_{a}$\n(where $e^{a}$ and $T^{a}$ are the elfbein and torsion, respectively),\nwhich might lead to inconsistencies~\\cite{Ede08}.\nOur own approach\nsuggests that\n$A_{3}$ should to be related to the CS fields through a much more complicated ansatz\n[see eq.~(\\ref{eq:A3})].\n\nRather than focusing on the dynamics of the TZ theory, in this paper we deal with its Lagrangian structure.\nThis proves to be a more tractable approach, since CS theories possess highly nonlinear dynamics and a complicated phase space structure~\\cite{Mis03}.\n\nWe first show that the TZ Lagrangian may be cast as a polynomial of degree~11 in $1\/l$, where $l$ is a length.\nWhen the CS five-index one-form $b^{abcde}$ vanishes, the leading term is just a cosmological constant,\nwhile the sub-leading term includes both torsion and a mass term for the fermions.\nOur main results follows:\nthe next term (proportional to $1\/l^{9}$) is essentially the CJS Lagrangian (see section~\\ref{sec:main}).\nThis discovery shows the existence of a previously unknown relation between the TZ and the CJS theories.\n\nThat this result may be new owes much to the great complexity of the TZ Lagrangian (see section~\\ref{sec:TZ}).\nWhile a compact closed formula exists [see eq.~(\\ref{eq:TZLag})],\na sufficiently explicit version would amount to roughly a thousand terms and has never been attempted.\n\nTo achieve our result we have first split the TZ Lagrangian into several meaningful terms\nby means of the CS subspace separation method introduced in Ref.~\\cite{Iza06a}.\nA careful dimensional analysis has then allowed us to extract only those terms proportional to the required powers of $l$.\nFor the sake of simplicity, we have restricted ourselves to the case when $b^{abcde} = 0$, which is roughly equivalent to setting $A_{3} = 0$.\nThe details of the calculation are summarized in~\\ref{sec:CJSTZ}.\n\nFurther comments and an outlook for future work are given in section~\\ref{sec:final}.\n\n\n\\section{A brief look at the TZ theory}\n\\label{sec:TZ}\n\nThe TZ theory is a pure gauge theory in eleven-dimen\\-sional spacetime\nwhose Lagrangian $\\mathcal{L}_{\\text{TZ}}$ is the CS eleven-form for the \\ensuremath{\\mathfrak{osp} \\left( 32 \\middle\\vert 1 \\right)}\\ superalgebra.\nA compact closed formula for $\\mathcal{L}_{\\text{TZ}}$ reads%\n\\footnote{Wedge product between differential forms is assumed throughout, e.g., $\\bm{A}^{2} = \\bm{A} \\wedge \\bm{A}$.}\n\\begin{equation}\n \\mathcal{L}_{\\text{TZ}} =\n 6 \\int_{0}^{1} dt\n \\left\\langle\n \\bm{A}\n \\left(\n t \\mathrm{d} \\bm{A} +\n t^{2} \\bm{A}^{2}\n \\right)^{5}\n \\right\\rangle,\n \\label{eq:TZLag}\n\\end{equation}\nwhere $\\bm{A}$ is an \\ensuremath{\\mathfrak{osp} \\left( 32 \\middle\\vert 1 \\right)}-valued one-form and\n$\\left\\langle \\cdots \\right\\rangle$\nstands for an \\ensuremath{\\mathfrak{osp} \\left( 32 \\middle\\vert 1 \\right)}-invariant supersymmetric polynomial of rank six.%\n\\footnote{This is essentially a higher-order version of the Killing metric, which may include a number of dimensionless coupling constants.}\n\nThe general properties of CS forms~\\cite{Che74,deAz95} ensure that $\\mathcal{L}_{\\text{TZ}}$\nchanges at most by a closed form under \\ensuremath{\\mathfrak{osp} \\left( 32 \\middle\\vert 1 \\right)}\\ gauge transformations, which close off-shell\nwithout introducing auxiliary fields.\n\nChoosing a convenient basis for \\ensuremath{\\mathfrak{osp} \\left( 32 \\middle\\vert 1 \\right)}\\ one may write $\\bm{A}$ as\n\\begin{equation}\n \\bm{A} =\n \\frac{1}{l} e^{a} \\bm{P}_{a} +\n \\frac{1}{2} \\omega^{ab} \\bm{J}_{ab} +\n \\frac{1}{5!l} b^{abcde} \\bm{Z}_{abcde} +\n \\frac{1}{\\sqrt{l}} \\bar{\\bm{Q}} \\psi,\n \\label{eq:A}\n\\end{equation}\nwhere\n$\\bm{P}_{a}$ are (AdS nonabelian) translations,\n$\\bm{J}_{ab}$ are Lorentz rotations,\n$\\bm{Z}_{abcde}$ are extra bosonic generators (not central charges---$\\bm{Z}_{abcde}$ behaves as a five-index, fully antisymmetric Lorentz tensor), and\n$\\bm{Q}$ is a fermionic generator.\nThis basis, with its associated commutation relations,\nallows the interpretation of\n$e^{a}$ as the orthonormal frame of one-forms%\n\\footnote{The spacetime metric $g_{\\mu \\nu}$ is not a fundamental field in the TZ theory,\nbeing derived from $e^{a}$ by means of the relation\n$g_{\\mu \\nu} = \\eta_{ab} \\nwse{e}{a}{\\mu} \\nwse{e}{b}{\\nu}$.}\n(the \\textit{elfbein} one-form) and\n$\\omega^{ab}$ as the spin connection,\nwhich shows that gravity (in its first-order formulation) is included in the TZ theory.\n\nDimensional consistency demands that we introduce a length parameter $l$ in eq.~(\\ref{eq:A}).\nThis length scale allows the one-form $\\bm{A}$ to remain dimensionless while giving the right dimensions to the component fields\n$e^{a}$, $\\omega^{ab}$, $b^{abcde}$ and $\\psi$.\n\nThe definition of $\\mathcal{L}_{\\text{TZ}}$\nis only complete after we specify the precise meaning of\n$\\left\\langle \\cdots \\right\\rangle$.\nPerhaps the simplest way to build this invariant polynomial\nis to use the trace of the supersymmetrized product of six generators in a suitable supermatrix representation.%\n\\footnote{This is not the most general choice, though, and more general invariant polynomials may prove useful. We return to this in section~\\ref{sec:final}.}\nOur choice for such supermatrix representation is as follows:\n\\begin{align}\n\\bm{P}_{a} & = \\left[\n\\begin{array}[c]{cc}\n\\frac{1}{2} \\nwse{\\left( \\Gamma_{a} \\right)}{\\alpha}{\\beta} & 0 \\\\\n0 & 0\n\\end{array}\n\\right], \\\\\n\\bm{J}_{ab} & =\\left[\n\\begin{array}[c]{cc}\n\\frac{1}{2} \\nwse{\\left( \\Gamma_{ab} \\right)}{\\alpha}{\\beta} & 0 \\\\\n0 & 0\n\\end{array}\n\\right], \\\\\n\\bm{Z}_{abcde} & =\\left[\n\\begin{array}[c]{cc}\n\\frac{1}{2} \\nwse{\\left( \\Gamma_{abcde} \\right)}{\\alpha}{\\beta} & 0 \\\\\n0 & 0\n\\end{array}\n\\right], \\\\\n\\bm{Q}^{\\rho} & = \\left[\n\\begin{array}[c]{cc}\n0 & C^{\\alpha \\rho} \\\\\n\\delta_{\\beta}^{\\rho} & 0\n\\end{array}\n\\right],\n\\end{align}\nwhere $\\Gamma_{a}$, $\\Gamma_{ab}$ and $\\Gamma_{abcde}$ are Dirac Gamma matrices in $d=11$,\nand $C_{\\alpha \\beta} = -C_{\\beta \\alpha}$ is the charge conjugation matrix.\n\nThe computation of\n$\\left\\langle \\cdots \\right\\rangle$\nis thus reduced to the calculation of traces of products of Gamma matrices,\nwhich is conceptually straightforward but may be quite challenging in practice.\nIn fact, the invariant polynomial so-defined has yet to be fully computed explicitly,\nwhich means that the detailed structure of the TZ theory is not well known.\n\n\n\\section{CJS Supergravity and the TZ Theory}\n\\label{sec:main}\n\nThe inclusion of the length parameter $l$ in eq.~(\\ref{eq:A}) and the algebraic structure of CS Lagrangians\n[see eq.~(\\ref{eq:TZLag})]\nimply that the TZ Lagrangian can be written as a polynomial of degree~11 in $1\/l$, namely\n\\begin{equation}\n\\mathcal{L}_{\\text{TZ}} =\n\\sum_{p=0}^{11}\nl^{-p}\n\\mathcal{L}^{\\left( p \\right)}.\n\\label{eq:LTZpol}\n\\end{equation}\nWe expect $l$ to be\na small quantity.\nThe most relevant terms in eq.~(\\ref{eq:LTZpol}) are then\n\\begin{equation}\n\\mathcal{L}_{\\text{TZ}} =\nl^{-11} \\mathcal{L}^{\\left( 11 \\right)} +\nl^{-10} \\mathcal{L}^{\\left( 10 \\right)} +\nl^{-9} \\mathcal{L}^{\\left( 9 \\right)} +\n\\cdots.\n\\label{eq:LTZ3}\n\\end{equation}\nWithout yet evaluating the invariant polynomial\n$\\left\\langle \\cdots \\right\\rangle$,\nthe first two terms are given by\n(see~\\ref{sec:CJSTZ})\n\\begin{align}\n\\mathcal{L}^{\\left( 11 \\right)} & =\n \\frac{6}{11}\n \\left\\langle\n \\bm{e}\n \\left(\n \\bm{e}^{2}\n \\right)^{5}\n \\right\\rangle, \\label{eq:L11sim} \\\\\n\\mathcal{L}^{\\left( 10 \\right)} & =\n 15\n \\left\\langle\n \\bm{\\psi}\n \\left[ \\bm{e} , \\bm{\\psi} \\right]\n \\left(\n \\bm{e}^{2}\n \\right)^{4}\n \\right\\rangle +\n 3\n \\left\\langle\n \\bm{eT}\n \\left( \\bm{e}^{2} \\right)^{4}\n \\right\\rangle, \\label{eq:L10sim}\n\\end{align}\nwhere, for simplicity, we have set $\\bm{b} = 0$.%\n\\footnote{Here we use the abbreviations\n$\\bm{e} = e^{a} \\bm{P}_{a}$,\n$\\bm{\\omega} = \\frac{1}{2} \\omega^{ab} \\bm{J}_{ab}$,\n$\\bm{b} = \\frac{1}{5!} b^{abcde} \\bm{Z}_{abcde}$, and\n$\\bm{\\psi} = \\bar{\\bm{Q}} \\psi$.\nAll powers of $l$ are displayed explicitly.}\nAs advertised, the leading term is a cosmological constant,\nwhile the sub-leading term includes both torsion and a mass term for the fermions.\n\nThe evaluation of the next term\n(the one proportional to $1\/l^{9}$)\nis significantly more involved,\nsince it requires the computation of many components of\n$\\left\\langle \\cdots \\right\\rangle$.\nWe have left all details concerning this calculation for~\\ref{sec:CJSTZ},\nwhere we also explain the notation and conventions used below.\nThe result reads\n\\begin{align}\n\\mathcal{L}^{\\left( 9 \\right)} & =\n- \\frac{1}{4 \\times 9! \\kappa_{2}^{2}} \\varepsilon_{a_{1} \\cdots a_{11}} R^{a_{1} a_{2}} e^{a_{3}} \\cdots e^{a_{11}}\n- \\frac{7i}{4} \\bar{\\psi} \\Gamma_{\\left( 8 \\right)} \\mathrm{D} \\psi\n- \\frac{1}{2} ab\n+ \\nonumber \\\\ &\n+ \\frac{i}{8} \\left( T^{a} - \\frac{i \\kappa_{1}^{2}}{4} \\bar{\\psi} \\Gamma^{a} \\psi \\right) e_{a} \\bar{\\psi} \\Gamma_{\\left( 6 \\right)} \\psi\n+ \\frac{7}{2^{4} \\times 5} a \\left. \\ast a \\right.\n- \\frac{1}{2^{10} \\times 3^{2}} \\frac{\\left( 4! \\right)^{5}}{6!} b \\left. \\ast b \\right.,\n\\label{eq:L9}\n\\end{align}\nwhere we have also set $\\bm{b} = 0$.\n\nCompare now eq.~(\\ref{eq:L9}) with the CJS Lagrangian~\\cite{Jul99},\n\\begin{align}\n\\mathcal{L}_{\\text{CJS}} & =\n- \\frac{1}{4 \\times 9! \\kappa^{2}} \\varepsilon_{a_{1} \\cdots a_{11}} R^{a_{1} a_{2}} e^{a_{3}} \\cdots e^{a_{11}}\n+ \\frac{i}{2} \\bar{\\psi} \\Gamma_{(8)} \\mathrm{D} \\psi\n- \\frac{1}{2} ab\n+ \\nonumber \\\\ &\n+ \\frac{i}{8} \\left( T^{a} - \\frac{i \\kappa^{2}}{4} \\bar{\\psi} \\Gamma^{a} \\psi \\right) e_{a} \\bar{\\psi} \\Gamma_{\\left( 6 \\right)} \\psi\n- \\frac{1}{2} \\left( a + F \\right) \\left. \\ast F \\right.,\n\\label{eq:LCJS}\n\\end{align}\nwhere we have set $A_{3} = 0$.\n\nThe similarity between eqs.~(\\ref{eq:L9}) and~(\\ref{eq:LCJS}) is, of course, striking.\nApart from some slightly different numerical coefficients,\nthe main difference is to be found among the last few terms,\nwhich are precisely those more sensitive to the truncation made by setting\n$\\bm{b} = 0$ and $A_{3} = 0$.\n\nIt could be argued that this similarity is just accidental,\nand that, given the highly nonlinear structure of the TZ Lagrangian,\nalmost any imaginable kind of term will appear in it.\nInterestingly enough, this is not the case:\nthere are many possible combinations which could in\nprinciple appear in the $l^{-9}$ sector of the TZ Lagrangian,\nand which do not show up in the CJS supergravity Lagrangian.\nNone of these appear.\nSome examples are\n\\begin{align}\n\\left[ EB \\right]^{abcde} & = \\nwse{E}{abcd}{f} B^{fe}, \\label{noshowfirst} \\\\\n\\left[ EBBB \\right]^{abcde} & = \\nwse{E}{ab}{fgh} B^{fc} B^{gd} B^{he}, \\\\\n\\left[ EBB^{2} \\right]^{abcde} & = \\nwse{E}{abc}{fg} B^{fd} \\nwse{B}{g}{h} B^{he}, \\\\\n\\left[ EB^{3} \\right]^{abcde} & = \\nwse{E}{abcd}{f} \\nwse{B}{f}{g} \\nwse{B}{g}{h} B^{he}, \\\\\n\\left[ BEB \\right]^{abc} & = B_{de} \\nwse{E}{deab}{f} B^{fc}, \\\\\n\\left[ BBEB \\right]^{a} & = B_{bc} B_{de} \\nwse{E}{bcde}{f} B^{fa}, \\label{noshowlast}\n\\end{align}\nwhere $B^{ab}$ is any antisymmetric two-index tensor constructed from the physical fields\n(e.g., $e^{a} e^{b}$, $R^{ab}$, etc.).\nSimilarly, $E^{abcde}$ stands for any antisymmetric five-index combination, for which,\neven with $b^{abcde}=0$, there are several possibilities:\n$\\bar{\\psi} \\Gamma^{abcde} \\psi$, $e^{a} e^{b} e^{c} e^{d} e^{e} $, etc.\n\n\n\n\\section{Discussion}\n\\label{sec:final}\n\nWe have shown that the TZ Lagrangian may be cast as a polynomial of degree 11 in $1\/l$,\nwhere $l$ is a (small) length parameter,\nand that the term proportional to $1\/l^{9}$\nis essentially identical to the CJS Lagrangian.\n\nOne significance of our result lies on its improbability:\nthe TZ theory and CJS supergravity emerge from entirely different backgrounds,\nyet both seem to be related at a deep level.\nWe would like to emphasize this:\nto the best of our knowledge,\nthere is a priori no reason whatsoever for the TZ theory to be related in this way to CJS supergravity,\nyet the relation appears.\n\nFor instance, the CJS Lagrangian includes the Hodge $\\ast$-operator,\nwhereas the TZ Lagrangian, being formulated with no a priori background metric, does not.\nThis sole fact could be na\\\"{\\i}vely argued to provide a reason for the absence of a relation between both theories.\nHowever, terms such as $a \\left. \\ast a \\right.$ do end up appearing in the TZ Lagrangian.\nThis can be traced back to the fact that a duality relation exists between different Dirac Gamma matrices in eleven-dimensional spacetime.\n\nA further interesting aspect deals with the \\ensuremath{\\mathfrak{osp} \\left( 32 \\middle\\vert 1 \\right)}\\ symmetry:\nour results suggest that the CJS Lagrangian fails to be invariant under \\ensuremath{\\mathfrak{osp} \\left( 32 \\middle\\vert 1 \\right)}\\\nbecause it lacks all the other terms that ensure the invariance of the TZ Lagrangian.%\n\\footnote{It has been known for some time, however, that CJS supergravity has on-shell symmetry under a one-parameter family of superalgebras~\\cite{DAu82,Ban04a,Ban04b}.}\n\nThe cosmological constant issue is seen here under a new light.\nWhile it is well known that CJS supergravity forbids a cosmological constant,\nthe TZ theory includes it on a different level, namely, in the term proportional to $1\/l^{11}$\nrather than in the one proportional to $1\/l^{9}$ (where the CJS terms reside).\nThe implications of this fact have yet to be fully assessed.\n\nThese observations lead us to propose the conjecture that CJS supergravity might actually be a truncation of a more general theory, such as the TZ theory.\nIn fact, the possibility that the CJS theory may be included within the TZ theory was already put forward by its authors~\\cite{Tro97}.\n\nOur result also highlights the usefulness of the CS subspace separation method~\\cite{Iza06a}\nfor manipulating CS Lagrangians in higher-dimensional spacetime.\nGiven the complexities involved, such powerful methods become essential tools\nto extract any meaningful information from CS theories (see~\\ref{sec:CJSTZ}).\n\nThere are several ways in which this work could be extended.\nThe inclusion of a nonzero $b^{abcde}$ is obviously one of them,\nbut the technical hurdles involved may be hard to overcome.\nSince only one-forms are allowed in CS theories,\nthe $A_{3}$ three-form must be first decomposed in terms of the TZ fields before a comparison is attempted.\nThe idea of decomposing $A_{3}$ in terms of one-forms has been already analyzed in the literature\n(see, e.g., Refs.~\\cite{Zan05,Ban04b}).\nOur construction has, however, an additional constraint;\n$A_{3}$ must be of order $l^{-3}$,\nand therefore, there are forbidden combinations.\nFor instance, the gravitino cannot be part of $A_{3}$,\nbecause it will lead to orders $l^{-1}$ and $l^{-2}$.\nDerivatives of the fields are also forbidden,\nsince it would lead at most to order $l^{-2}$ and order $l^{-3}$ would never be reached.\nTherefore, the identification of $A_{3}$ with a term such as $e^{a} T_{a}$ \\cite{Banh01}\nis ruled out from the outset.\nA moment's thought shows that the only consistent way of decomposing $A_{3}$\nin terms of the TZ fields is through the ansatz\n\\begin{align}\nA_{3} & =\nC_{abc} e^{a} e^{b} e^{c} +\n\\frac{1}{5!} C_{abc_{1} \\cdots c_{5}} e^{a} e^{b} b^{c_{1} \\cdots c_{5}}\n+ \\frac{1}{\\left( 5! \\right)^{2}} C_{ab_{1} \\cdots b_{5} c_{1} \\cdots c_{5}} e^{a} b^{b_{1} \\cdots b_{5}} b^{c_{1} \\cdots c_{5}} +\n\\nonumber \\\\ & +\n\\frac{1}{\\left( 5! \\right)^{3}} C_{a_{1} \\cdots a_{5} b_{1} \\cdots b_{5} c_{1} \\cdots c_{5}} b^{a_{1} \\cdots a_{5}} b^{b_{1} \\cdots b_{5}} b^{c_{1} \\cdots c_{5}},\n\\label{eq:A3}\n\\end{align}\nwhere the various $C$-tensors are fully anti-symmetric zero-forms to be determined.\n\nAs mentioned in section~\\ref{sec:TZ}, the invariant polynomial we have used is but the simplest choice.\nMore general invariant polynomials may prove useful for eliminating the slight differences that still remain between the CJS Lagrangian and the $1\/l^{9}$ sector of the TZ theory.\n\nAnother interesting issue is the exploration of the solution space of the TZ theory.\nWhile the full theory is difficult to analyze due to the nonlinearities of the Lagrangian~\\cite{Banh95,Banh96b,Cha99,Mis03},\nthe truncated version (at order $1\/l^{9}$) should prove to be more tractable.\n\n\n\\section*{Acknowledgments}\n\nThe authors wish to thank\nJos\u00e9~A.\\ de~Azc\u00e1rraga,\nAlfredo P\u00e9rez,\nPatricio Salgado and\nJorge Zanelli\nfor many fruitful discussions regarding the topics covered in the present work.\nWe are also indebted to\nRicardo Troncoso\nfor helpful comments on the first draft of this manuscript.\nF.~I.\\ extends his thanks to\nJos\u00e9~A.\\ de~Azc\u00e1rraga\nfor his kind hospitality at the Department of Theoretical Physics of the Universitat de Val\\`{e}ncia, and to\nMicho {\\DJ}ur{\\dj}evich\\\nfor his kind hospitality at the Institute of Mathematics of the Universidad Nacional Aut\u00f3noma de M\u00e9xico, where part of the present work was done.\nE.~R.\\ is grateful to\nJorge Zanelli\nfor his kind hospitality at the Centro de Estudios Cient\u00edficos (CECS), Valdivia, where part of this work was done.\nThe authors also thank\nKasper Peeters\nfor his computer algebra system ``Cadabra'' \\cite{Pee06,Pee07}, which was of invaluable help when performing lengthy Dirac-matrix algebra computations, and\nRicardo Ram\u00edrez,\nfor helping to create the necessary code to evaluate some of the components of the invariant polynomial.\nF.~I.\\ and E.~R.\\ were supported by the\nNational Commission for Scientific and Technological Research, Chile, through Fondecyt research grants 11080200 and 11080156, respectively.\n\n\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\t \nThe majority of stars do not form in isolation but in areas of clustered star formation \\citep[e.g.][]{lada2003}. In some cases, the densest areas of these large regions result in gravitationally bound stellar systems, commonly referred to as star clusters. These bound systems can survive for hundreds of Myr. To distinguish them from ancient stellar objects like the globular clusters (GCs), we refer to them as young stellar clusters (YSCs). They usually populate star-forming galaxies in the local universe (e.g. \\citealp{larsen2006b}) and their physical properties (ages, masses) can in principle be used to determined star formation histories (SFHs) of the hosting galaxies \\citep[e.g.][]{miller1997,goudfrooij2004,konstantopoulos2009,glatt2010}.\n\nOver the past 20 years, studies of the distributions of YSC luminosities and masses in local galaxies have shown that they are well-described by a power law function of the form $Ndm \\propto M^\\alpha dm$, with a slope $\\alpha \\sim-2$, observed both for low-mass clusters in the Milky Way \\citep{piskunov2006}, in the Magellanic Clouds \\citep{baumgardt2013,degrijs2006} and in M31 \\citep{fouesneau2014}, and for sources up to masses of $\\sim10^5-10^6$ M$_\\odot$\\ in nearby spirals and starburst galaxies \\citep{chandar2010,konstantopoulos2013,whitmore2010}. This result is expected if star formation happens in a hierarchical manner, dominated by interstellar medium (ISM) turbulence, and the clusters occupy the densest regions (e.g. \\citealp{elmegreen2006}; see also \\citealp{elmegreen2010b} for a review).\n\nDespite observational and theoretical progress over the past few decades, many questions concerning the properties of YSC populations remain open. Among these: is cluster formation only driven in space and time by size--of--sample effects \\citep[e.g.][]{hunter03}, with an increasing number of clusters found in galaxies with higher star formation rate (SFR)? Will the galactic environment (ISM conditions, gas fraction, galaxy type) where clusters form leave an imprint on the final properties of the YSC populations? When we look at YSC populations in local spirals \\citep[e.g.][]{larsen2004}, merger systems \\citep[e.g.][]{whitmore1999} and dwarf galaxies \\citep[e.g.][]{billett2002} it is challenging to discern the role played by statistical sampling \\citep[e.g.][]{fumagalli2011} and environment.\n\nEven the exact shape of the mass function is still debated, in particular concerning its high mass end. Some early studies \\citep[e.g.][]{larsen2006} have pointed out the dearth of massive YSCs if a single power law fit of slope $-2$ describes the upper-part of the YSC mass function. \\cite{gieles2006} have proposed a Schechter function as a better description of the YSC mass function in local galaxies, due to a mass truncation at a characteristic mass above which the likelihood of forming massive clusters goes rapidly to zero. \n\n\nIn order to be able to characterise how star clusters form and evolve, \nit is important to study a statistically meaningful sample. The Legacy Extragalactic UV Survey (LEGUS) is a Cycle 21 HST Treasury program which observed 50 nearby galaxies from the UV to NIR bands, with the goal of deriving high quality star cluster catalogues, and, more in general, of studying star formation at intermediate scales, linking the smallest (stellar) scales to the larger (galactic) ones \\citep[see][]{legus1}. \nIn general, the large number of galaxies and galaxy properties available in LEGUS will enable us to statistically study YSC populations over a wide range of galactic environments \\citep{legus2}. \n\nAmong the most interesting galaxies in the LEGUS catalogue is NGC 5194 (also known as M51a or the Whirlpool Galaxy), because of its proximity and the number of star clusters that it hosts.\nIt is a spiral galaxy, catalogued as SAbc\\footnote{According to the Nasa Extra-galactic Database (NED)}, almost face-on (inclination angle $i\\approx22^\\circ$, \\citealp{colombo2014b}) at a distance of 7.66 Mpc \\citep{distancem51}. M51a is interacting with the (smaller) companion galaxy NGC 5195 and it is probably this interaction that is the cause of a marked spiral geometry and a high star formation process (a SFR value of 2.9 M$_\\odot$\/yr is derived from published total fluxes in the far-$UV$ and 24$\\mu$m, combined using the recipe by \\citealp{hao2011}) sustained over time \\citep[e.g.][]{dobbs2010_m51}. The two galaxies together form the M51 system. In the remainder of this paper we will use the name M51 mainly referring to the main spiral galaxy M51a.\nThis galaxy hosts numerous star formation complexes \\citep{bastian2005b}, HII regions \\citep{thilker2000,lee2011} and YSCs and it has been a benchmark in the study of extragalactic star and cluster formation. \n\nHigh-brightness blue sources in M51 have been studied already by \\citet{georgiev1990}.\nIn more recent years, broadband and narrowband imaging with the Hubble Space Telescope (HST) Wide-Field Planetary Camera 2 (WFPC2) in various bands from $UV$ to $NIR$ were used for initial studies of the cluster population in small parts of the galaxy \\citep{bik2003,bastian2005,gieles2005,lee2005}. Later optical observations with the higher resolution and more sensitive ACS camera were obtained in the $BVI$ bands and covered uniformly the entire galaxy allowing to extend the investigation of the YSC population to the whole galaxy \\citep{scheepmaker2007,hwang_lee2008, chandar2011}. More recently, the coverage by the WFPC2 F336W filter ($U$-band) has been expanded, with 5 more pointings, along with H$\\alpha$\\, data, allowing improved age determination for a significant fraction of the cluster population\\footnote{The $U$-band filter (or bluer filters) is fundamental to break the age-extinction degeneracy when SEDs are compared to stellar population synthesis models, see \\citet{anders2004}} \\citep{scheepmaker2009,hwang_lee2010,chandar2011,chandar16}.\n\nAll this effort led to the consensus that the star cluster population in M51 can be described by a standard mass distribution, i.e. a simple power law with slope -2. However, whether the single-power law function is also a good representation of the upper-mass end of the cluster mass function, in terms of the eventual presence of a truncation at high masses, is still under debate (compare e.g. \\citealp{gieles2006} and \\citealp{chandar2011}). \nThe analyses of the cluster mass function evolving in time, and, more in general, of the cluster number densities evolving with time, reach different conclusions on the disruption properties of the clusters in M51. Some studies observe a mass function evolution consistent with a disruption time dependent on the mass of the clusters \\citep[e.g. the mass-dependent disruption -MDD- model by][]{gieles2009}, while in others the study of the mass function (MF) evolution seems to exclude this model, and to favour a constant disruption time of clusters \\citep[e.g. mass independent disruption -MID- model by][]{chandar16}. \n\nThe interaction of M51 has been studied using simulations in order to describe the current geometrical and dynamical properties of the star formation \\citep{salo2000,dobbs2010_m51}. \nCluster properties have then been compared with the expectations based on simulations in order to test the models for the formation of the spiral structure (e.g. \\citealp{chandar2011} ruled out the possibility of self-gravity as the cause of the generation of the spiral structure). \n\nStar formation in M51 has also been studied from the point of view of molecular gas via radio observations (\\citealp{schuster2007,koda2009,koda2011} and \\citealp{schinnerer2010,schinnerer13} among the most recent). High resolution interferometric data have been used to study in detail the properties of giant molecular clouds (GMCs) \\citep{koda2012,colombo2014a}.\nThe possibility of studying the galaxy at high resolution at different wavelengths allows studying star formation at different ages, in particular to compare the properties of the progenitors (GMCs) and the final products (stars and star clusters). \n\nOne of the goals of the present work is to conduct a statistically driven study of the YSC population of M51 using the new data and cluster catalogue produced by the LEGUS team. The new LEGUS dataset of M51 provides 5 new pointings in the $NUV$ (F275W and F336W) with the Wide Field Camera 3 (WFC3). The improved spatial resolution of the WFC3 and sensitivity in the $NUV$ give a better leverage on the physical determinations of the YSC properties \\citep{legus1}.\nIn order to compare our new catalogue with previously published works we investigate, in this paper, YSC mass and luminosity functions for the whole galaxy. With the help of simulated Monte Carlo cluster populations we build a comprehensive picture of the cluster formation and evolution in the galaxy as a whole. In a forthcoming paper (Messa et al., in prep, hereafter Paper II) we test whether YSC properties change across the galaxy as a function of star formation rate (SFR) density ($\\Sigma_{\\textrm{SFR}}$) and gas surface density. These results can shed light on a possible environmental dependences in the properties of the cluster population and whether studies of YSC populations can be used to characterise the galactic environment.\n\nThe paper is divided as follows: a short description of the data is given in Section~\\ref{sec:data} and the steps necessary to produce the final cluster catalogue are described in Section~\\ref{sec:ccp}. In Section~\\ref{sec:global} the global properties of the sample (luminosity, mass and age functions) are studied, while in Section~\\ref{sec:montecarlo} the same properties are analysed using simulated Monte Carlo populations. The fraction of star formation happening in a clustered fashion is studied in Section~\\ref{sec:cfe}.\nFinally, the conclusions are summarised in Section~\\ref{sec:conclusions}. \n\n\n\\section {Data}\n\\label{sec:data}\nA detailed description of the LEGUS general dataset and the standard data reduction used for LEGUS imaging is given in \\citet{legus1} and we refer the reader to that paper for details on the data reduction steps.\n\nHere we summarize the properties of the data used in this study. The M51 system (NGC 5194 and NGC 5195) spans $\\sim7 \\times 10$ arcmin on the sky at optical wavelengths (at an assumed distance of 7.66 Mpc, from \\citealp{distancem51}) and several pointings are therefore necessary to cover their entire angular size. \nThe LEGUS dataset includes multi-band data spanning the wavelength range from near-$UV$ to near-$IR$; data for M51 cover the $UV$ (F275W), $U$ (F336W), $B$ (F435W), $V$ (F555W) and $I$ (F814W) bands. Even though no conversion is applied to the Cousins-Johnson filter system, we keep the same nomenclature, due to the similarity of the central wavelength between that system and our data. Concerning the $B$, $V$ and $I$ filters, ACS WFC archival data\navailable from the Mikulski Archive for Space Telescopes (MAST) have been re-processed. The data in these bands include 6 pointings that cover the entire galaxy and the companion galaxy NGC 5195 (GO-10452, PI: S. Beckwith).\n\nWithin the LEGUS project, the coverage has been extended to the $U$ and $UV$ bands. The new $UV$ data consist of 4 pointings covering the arms and outskirts of the galaxy combined with a deep central exposure (GO-13340, PI: S. Van Dyk) covering the nuclear region of the galaxy. Exposure times for all filters are summarized in Tab. \\ref{tab:data} while the footprints of the pointings are illustrated in Figure \\ref{fig:footprints}.\n\n\\begin{table}\n\\centering\n\\caption{Exposure times for the different filters and number of pointings (the exposure times refer to each single pointing). As can be noted also in Fig.~\\ref{fig:footprints} the ACS data cover the entire galaxy with 6 pointings, while for the $UV\/U$ band the observations consist of 5 pointings only.}\n\\begin{tabular}{clccl}\n\\hline\nInstr. & Filter & Expt. & \\# & Project Nr. \\& PI \\\\\n\\hline\n\\hline\nWFC3\t& F275W($UV$)& 2500 s\t& 4\t& GO-13364\\ \\ \t D. Calzetti \\\\\n\\ \t\t & \\\t\t\t& 7147 s\t& 1\t& GO-13340\\ \\ \t S. Van Dyk \\\\\nWFC3\t& F336W($U$)\t& 2400 s\t& 4\t& GO-13364\\ \\\t D. Calzetti \\\\\n\\\t\t & \\\t\t\t& 4360 s\t& 1\t& GO-13340\\ \\ \t S. Van Dyk \\\\\nACS \t\t& F435W($B$)\t& 2720 s\t& 6\t& GO-10452\\ \\ S. Beckwith \\\\\nACS\t\t& F555W($V$)\t& 1360 s\t& 6\t& GO-10452\\ \\ \t S. Beckwith \\\\\nACS\t\t& F814W($I$)\t& 1360 s\t& 6\t& GO-10452\\ \\ \t S. Beckwith \\\\\n\\hline\n\\end{tabular}\n\\label{tab:data}\n\\end{table}\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=0.4\\textwidth]{01_ngc5194_footprints.png}\n\\caption{UVIS (red) and ACS (orange) footprints on a DSS image of the NGC5194 and NGC5195 system. The UVIS (white) footprint corresponds to proposal 13340 (PI: S. Van Dyk). See also Tab. \\ref{tab:data} for more info on the observations.}\n\\label{fig:footprints}\n\\end{figure}\n\n\n\\section{Cluster Catalogue production}\n\\label{sec:ccp}\n\\subsection{Cluster Extraction}\n\\label{sec:cluster_extraction}\nIn order to produce a cluster catalogue of the M51 galaxy we follow the procedures described in \\cite{legus2}, where a detailed description of the standard reduction steps can be found. Hereafter we describe these steps along with the specific parameters used for the M51 dataset. The catalog\nproduction is divided into 2 main parts, the cluster extraction and the cluster classification.\n\nThe cluster extraction is executed through a semi-automatic custom pipeline available inside the LEGUS collaboration. As the first step we extracted the source position of the cluster candidates in the $V$ band (used as reference frame in our analysis) with \\texttt{SExtractor} \\citep{sextractor}. The parameters of \\texttt{SExtractor} were chosen to extract sources with at least a 10$\\sigma$ detection in a minimum of 10 contiguous pixels.\nIn the same band, we measured the concentration index (CI) on each of the extracted sources.\nWe use the definition for the CI as the magnitude difference between the fluxes in circular regions of 1 and 3 pixels radius, centred on the source position. It measures how much the light is concentrated in the centre of the source and can also be used as also a tracer of the cluster size \\citep[see][]{ryon2017}. \n\n\\begin{figure*}\n\\centering\n\\subfigure{\\includegraphics[width=0.9\\columnwidth]{02a_plot_histogram_ci_fullcat.pdf}}\n\\hspace{1mm}\n\\subfigure{\\includegraphics[width=0.9\\columnwidth]{02b_plot_histogram_ci.pdf}}\n\\subfigure{ \\includegraphics[width=0.18\\textwidth]{02c_plot_apcor_275_f275w.png}}\n\\subfigure{ \\includegraphics[width=0.18\\textwidth]{02d_plot_apcor_336_f336w.png}}\n\\subfigure{ \\includegraphics[width=0.18\\textwidth]{02e_plot_apcor_435_f435w.png}}\n\\subfigure{ \\includegraphics[width=0.18\\textwidth]{02f_plot_apcor_f555w.png}}\n\\subfigure{ \\includegraphics[width=0.18\\textwidth]{02g_plot_apcor_814_f814w.png}}\n\\caption{\\textit{Top left:} Distribution of the CI values for all the sources extracted with \\texttt{SExtractor}. The red solid line indicates the value of 1.35 used to select the cluster catalogue. \\textit{Top right:} CI distribution for visually selected stars (in red) and clusters (in blue). \\textit{Bottom:} Distributions of the aperture corrections given by visually selected clusters. Black dashed lines set the interval within which values are considered to calculate the average value (blue bin). \n}\n\\label{fig:histo_CI}\n\\end{figure*}\nThe distribution of the CI values for the extracted sources looks like a continuous distribution, peaked around a value of 1.3, as Fig.~\\ref{fig:histo_CI} (top left) shows, but is in fact the sum of two sub-distributions, one for stars and one for clusters. To understand how the distributions of CI values changes between stars and clusters we select via visual inspection a sample of stars and clusters that are used as training samples for our analysis. In Fig.~\\ref{fig:histo_CI}, top right, we show the CI distributions of stars and clusters. It becomes clear that the distributions of CIs in the two cases are different. Stars, being point sources, have CI values that do not exceed values of $\\sim$1.4, while clusters have on average higher CI values and a broader distribution. The distributions also suggest that considering sources with a CI bigger than 1.35 strongly decreases the chances of erroneously including stars in the catalogue, thus facilitating the selection of most of the clusters. Following the CI versus effective radius relation showed in Fig. 4 of \\citet{legus2}, we estimate that a CI cut at 1.35 mag corresponds to a cluster effective radius of 0.5 pc. Because the distribution star cluster radii peak at $\\sim$3 pc \\citep{ryon2017}, placing a CI cut at 1.35 mag will not negatively impact the recovery of clusters in this system.\n\nAperture photometry was performed on the CI-selected sample, using fixed apertures of 4 pixels radius, and local sky is subtracted from an annulus with 7 pixels (px) interior radius and 1 px width. \nA fixed aperture correction was estimated using the photometric data of the visually selected sample of clusters. The sample was adjusted in each filter in order to consider only isolated bright clusters with well defined PSF wings. The number of visually selected sources used in each filter are listed in Tab~\\ref{tab:corrections}. During the visual selection, sources were chosen to span different cluster sizes and to also include compact clusters. In this way the resulting aperture correction is not biased towards the very large clusters which are more easily detectable. For each source the single aperture correction was calculated subtracting the standard photometry (aperture: 4 px and sky at 7 px) with the total photometry within a 20 px radius (with a 1 px wide sky annulus at a radius of 21 px). The final correction in each filter was calculated taking the average of the values within an allowed range of values. The single aperture correction values of the selected sample along with the final mean value in each filter are plotted in Fig.~\\ref{fig:histo_CI} (bottom) and the values are also reported in Tab.~\\ref{tab:corrections}. The standard deviation ($\\sigma_{\\lambda}$) is added in quadrature to the photometric error of each source. \n\\begin{table}\n\\centering\n\\caption{Corrections applied to the photometry of all the sources in different filters. With reddening we refer to the Galactic extinction in magnitudes (in each filter). The uncertainty on the average aperture correction has been used to estimate the final error on the magnitude.}\\begin{tabular}{ccccc}\n\\hline\nFilter\t\t& Reddening\t& \\# Clusters \t& Avg ap.corr. \t& $\\sigma_{\\textrm{ap.corr.}}$\t\\\\\n\\\t\t& [mag]\t\t& \t\t\t& [mag]\t \t& [mag]\t\\\\\n\\hline\n\\hline\nF275W\t& 0.192\t\t& 36\t\t\t\t& -0.628 \t\t& 0.034 \\\\\nF336W\t& 0.156\t\t& 66\t\t\t\t& -0.668 \t\t& 0.030 \\\\\nF435W\t& 0.127\t\t& 62\t\t\t\t& -0.665 \t\t& 0.026 \\\\\nF555W\t& 0.098\t\t& 71\t\t\t\t& -0.663\t\t& 0.023 \\\\\nF814W\t& 0.054\t\t& 56\t\t\t\t& -0.830\t\t& 0.031 \\\\\n\\hline\n\\end{tabular}\n\\label{tab:corrections}\n\\end{table}\n\nA final cut was made excluding sources which are not detected in at least 2 contiguous bands (the reference $V$ band and either the $B$ or $I$ band) with a photometric error smaller than 0.3 mag. The positions of the 30176 sources satisfying the CI cut of 1.35 mag and this last selection criterion are collected, along with their photometric data, in a catalogue named \\textit{``automatic\\_catalog\\_avgapcor\\_ngc5194.tab''}, following the LEGUS naming convention. Note that, being automatically-selected, this catalogue probably includes contaminant sources (e.g. foreground stars, background galaxies, stars in the field of M51).\n\nTo remove the contamination of non-clusters in the automatic catalog, we created a high-fidelity sub-catalog including all sources detected in at least 4 bands with a photometric error below 0.30 mag and having an absolute $V$ band magnitude brighter than -6 mag. Selecting only bright sources reduces considerably the number of stars in the catalogue, while the constraint on the number of detected bands allows a reliable process for the SED fitting analysis (see Section~\\ref{sec:sed}). Note that, differently from the standard LEGUS procedure, we applied the -6 mag cut to the \naverage--aperture--corrected magnitudes and not to the CI--based--corrected ones (see \\citealp{legus2} for a description of the CI--based correction).\nThis choice is motivated mainly by the use of the average--aperture--corrected catalogue as the reference one: when testing the completeness level of our catalogue (see Section~\\ref{sec:completeness}) we noticed that applying the cut on the average--aperture--corrected magnitudes improves the completeness, being more conservative (i.e. allowing the inclusion of more sources). This high-fidelity catalogue counts 10925 sources, which have been all morphologically classified (see Section~\\ref{sec:classification}).\n\n\\subsection{Morphological classification of the cluster candidates, human versus machine learning inspection}\n\\label{sec:classification}\n\nSources in the high fidelity catalogue were visually inspected, in order to morphologically classify the cluster candidates and exclude stars and interlopers that passed the automatic selection.\nLike for the other galaxies of the LEGUS sample, visually inspected sources were divided into 4 classes, described and illustrated in \\citet{legus2}. Briefly, class 1 contains compact (but more extended than stars) and spherically symmetric sources while class 2 contains similarly compact sources but with a less symmetric light distribution. Both these classes include sources with a uniform colour. Class 3 sources show multi-peaked profiles with underlying diffuse wings, which can trace the presence of (small and compact) associations of stars. Sources in class 3 can have colour gradients. Contaminants like single stars, or multiple stars that lie on nearby pixels even if not part of a single structure, and background galaxies are all stored in class 4 and excluded from the study of the cluster population of the galaxy.\n\n\\begin{table}\n\\centering\n\\caption{Number of sources in each class for the human-classified (2nd column) and machine-learning classified (3rd column) sources. In parentheses are the percentage of sources over the number of the visually inspected sources. The fourth column lists the percentage of the sources in each class for which the ML has assigned the same class as the humans. In the fifth column the number of sources per class classified by the machine learning algorithm is given. In brackets we include the recovered fraction with respect to the total number of sources in the catalogue (i.e. 10925 sources).}\n\\begin{tabular}{cllll}\n\\hline\nClass \t& Human \t\t\t& ML\t\t\t\t&HvsML\t\t& ML\t (tot cat.) \\\\\n\\hline\n\\hline\ntot\t\t& 2487\t\t\t& 2487\t\t\t&\\ \t\t\t& 10925\t\t\\\\\n\\hline\n1\t\t& 360 ($14.5\\%$) \t& 377 ($15.2\\%$)\t& $95.3\\%\t$\t& 1324 ($12.1\\%$)\\\\\n2\t\t& 500 ($20.1\\%$) \t& 516 ($20.7\\%$)\t& $93.4\\%$\t& 1665 ($15.2\\%$)\\\\\n3\t\t& 365 ($14.7\\%$) \t& 338 ($13.6\\%$)\t& $92.1\\%$\t& 385 ($3.5\\%$)\\\\\n4\t\t& 1262 ($50.7\\%$)\t& 1256 ($50.5\\%$)\t& $95.6\\%$\t& 7551 ($69.1\\%$)\\\\\n\\hline\n\\end{tabular}\n\\label{tab:classification}\n\\end{table}\nDue to the large number of sources entering the automatic catalogue we have implemented the use of a machine learning (ML) optimised classification. We have visually inspected only $\\sim$1\/4 of the 10925 sources, located in different regions of the galaxy. This visually inspected subsample has been used as a training set for the ML algorithm to classify the entire catalogue (details of the algorithm are discussed in a forthcoming paper by Grasha et al., in prep). \nThe ML code is run on the entire sample of 10925 sources, including the already humanly--classified ones. In this way we can use the sources having a classification with both methods to estimate the goodness of the ML classification for M51. Table~\\ref{tab:classification} gives the number of sources classified in each class by human and ML, as well as the comparison between the two classification (forth column). We recover a 95\\% of agreement between the two different classifiers, within the areas used as training sets. To assess the goodness of the ML classification on the entire sample, we list in Table~\\ref{tab:classification} between brackets the relative fraction of each class with respect to the total number of sources classified with different methods. We observe that the relative number of class 1 and 2 sources with respect to the total number of sources (10925) classified by the ML approach is only a few per cent smaller than the fraction obtained with the control sample (2487 sources). However, there is a striking difference in the recovery fractions of class 3 and 4 sources. When considering the entire catalogue, the relative number of class 3 objects is much smaller (and on the contrary the one of class 4 is significantly more numerous). We consider very unlikely that there are so few class 3 objects in the total sample. So far the ML algorithm fails in recognising the most variate class of our classification scheme that contains sources with irregular morphology, multi-peaked, and some degree of colour gradient. From the absolute numbers of sources per class it is easy to conclude that the ML code is able to reclassify correctly almost all the class 3 objects given as a training sample, but is unable to recognize new class 3 sources, considering many of them as class 4 objects.\nFuture improvements for the classification will be produced by the use of different ML recognition algorithms. For our current analysis we will focus on the properties of class 1 and 2 cluster candidates and exclude class 3 sources, as explained in Sec \\ref{sec:finalcat}.\n\n\nWe can summarise the photometric properties of the M51 cluster population using a two colour diagram (Fig~\\ref{fig:ccd}).\n\\begin{figure}\n\\centering\n\\includegraphics[width=\\columnwidth]{03_M51_CCD_optical_bestm_cl12.png}\n\\caption{Colour--colour diagram with $V-I$ on the x axis and $U-B$ on the y axis. Black points are all the sources in the catalogue that passed the automatic selection, while the orange shaded area gives the density of only the class 1 ad 2 sources. The SSP evolutionary track from Padova-AGB models is overplotted. The ages covered by the track goes from 1 Myr to 14 Gyr. The arrow indicates how an object would move if corrected for a reddening E($B-V$)=0.2. The cross at the top-right shows the average error in colours. }\n\\label{fig:ccd}\n\\end{figure}\nContours based on number densities of clusters show the regions occupied by the class 1 and class 2 M51 cluster population with respect to the location of the 10925 sources included by the automatic selection. A simple stellar population (SSP) track showing the cluster colour evolution as a function of age is also included. Sources are mainly situated along the tracks, implying the high quality of our morphological classification. Extinction spreads the observed colours towards the right side of the evolutionary tracks. Correcting for extinction, in the direction indicated by the black arrow, would move the sources back on the track, at the position corresponding to the best--fitted age.\nThe diagram shows a broad peak between 50 and 100 Myr. The contours are quite shallow towards younger ages and show a pronounced decline around $\\sim1$ Gyr suggesting that most of the sources detected are younger than 1 Gyr. We use SED fitting analysis to derive cluster physical properties (see Section~\\ref{sec:sed}), including the age distribution (Fig~\\ref{fig:agemass}b). \n\n\n\n\\subsection{Completeness}\n\\label{sec:completeness}\n\nTo investigate the completeness limit of the final catalogue, we use the custom pipeline available within the LEGUS collaboration as described in \\citet{legus2}.\nThe pipeline follows closely the selection criteria adopted to produce the final automatic catalogues.\nFor each filter we produce frames containing simulated clusters of different luminosities and sizes which are added to the original scientific frames. Effective radii ($R_\\textrm{eff}$) between 1 and 5 pc have been used, as studies of cluster sizes in similar galaxies suggest that most of the sources fall in this range \\citep{ryon2015,ryon2017}. The synthetic clusters span an apparent magnitude range between 19 and 26 mag and are created using the software \\texttt{BAOlab}, freely available online\\footnote{\\hyperref[http:\/\/baolab.astroduo.org\/]{http:\/\/baolab.astroduo.org\/}}. All clusters are simulated as symmetric sources with a MOFFAT15 profile \\citep[see][]{baolab} considered a standard assumption for the YSC light profiles \\citep{eff1987}.\nCluster extraction (via \\texttt{SExtractor}) and photometry are repeated on the resulting coadded scientific and synthetic frames. \nA signal of 10$\\sigma$ in at least 10 contiguous pixel is requested to extract sources in the $B$, $V$, and $I$ bands, while a value of 5$\\sigma$ over at least 10 contiguous pixels is chosen for the $UV$ and $U$ bands.\nThe recovery rate of sources as a function of luminosity yields the completeness.\nThe magnitude limits above which the relative number of the recovered sources falls below 90\\% level is summarized in Tab.~\\ref{tab:completeness} for each filter. \n\\begin{table}\n\\centering\n\\caption{Completeness limits in each filter. The second and third column give the 90\\% completeness limit calculated with the completeness test code (described in the text). The completeness was computed in the disk (area outside a 35'' radius, second column) and in the central region (area inside 35'' radius, third column). The last column gives the peak magnitude of the luminosity distribution, as plotted in Fig~\\ref{fig:checkmaglimits}.}\n\\begin{tabular}{lccc}\n\\hline\nFilter\t\t& compl. (disk) & compl. (centre)\t& Lum$_\\textrm{peak}$\t\\\\\n\\\t\t& [mag] \t\t& [mag]\t\t\t& [mag]\t\\\\\n\\hline\n\\hline\nF275W\t& 22.17 \t& 21.75 \t& 21.75\t\\\\\nF336W\t& 22.75\t& 21.79\t& 22.00\t\\\\\nF435W\t& 24.17 \t& 22.65 \t& 23.25\t\\\\\nF555W\t& 23.70\t& 22.31\t& 23.25\t\\\\\nF814W\t& 22.70\t& 21.61\t& 22.25\t\\\\\n\\hline\n\\end{tabular}\n\\label{tab:completeness}\n\\end{table}\n\nThe completeness test code only gives a completeness limit estimated for each filter independently. These values should be considered as completeness limits resulting from the depth of the data. However, our cluster candidate catalogue is the result of several selection criteria that cross correlate the detection of sources among the 5 LEGUS bands. This effect can be visualized in Fig~\\ref{fig:checkmaglimits}, where we show the recovered luminosity distributions of sources at different stages of the data reduction. The requirement of detecting the sources in 4 filters with a photometric error smaller than 0.30 mag diminishes the number of recovered objects, mainly due to the smaller area covered by UVIS compared to ACS.\nThe cut at $M_V = -6$ mag not only excludes the sources which are faint in $V$ band, but also modifies the luminosity distributions in the other filters (see differences between red solid line distributions and blue dashed ones). Both these conditions affect the resulting catalogue and modify the shape of the luminosity distributions in each filter in a complicated way at the faint limits, therefore modifying also the completeness limits with respect to our approach of treating each filter independently. \nWe use the observed luminosity distributions to understand how the completeness limits change as a function of waveband and adopted selection criteria. \nWe observe in each filter an increase going from bright to low luminosities and we know that incompleteness starts to affect the catalogue where we see the luminosity distribution turning over (see e.g. \\citealp{larsen2002}). \n\\begin{figure*}\n \\centering\n \\includegraphics[width=\\textwidth]{04_plot_mag_lim_a.pdf}\n \\caption{Luminosity distributions showing how the completeness changes when different cuts are applied in the selection process. Each distribution represents the catalogue at a different stage: the black dotted line shows the distribution of all the sources with photometric errors smaller than 0.3 mag, the blue dashed line represents the sources left after requiring that they have a detection in at least 4 filters and the red solid line is the distribution of the sources after the -6 mag cut in the $V$ band. Vertical lines are plotted corresponding to the peak of the distribution. The peak values after the magnitude cut is applied in $V$ band are listed in Table~\\ref{tab:completeness}. The sources are grouped in bins of 0.25 mag width. The black (orange) arrows show the value of the 90\\% completeness limit evaluated with the completeness test code in the disk (centre) of the galaxy (see Tab.~\\ref{tab:completeness}).\n }\n\\label{fig:checkmaglimits}\n\\end{figure*}\n\nWe draw the following conclusions from the analysis conducted in Fig~\\ref{fig:checkmaglimits}. First, the $-6$ mag cut in the $V$ band strongly reduces the number of selected sources in all the bands. At the distance of M51, this brightness corresponds to an apparent magnitude in $V$ of 23.4 mag, thus, it is brighter than the 90\\% completeness limit recovered in the $V$ band (23.84 mag, see Tab.~\\ref{tab:completeness}). Secondly, the 90\\% completeness limits fall rightwards of the peak in the luminosity distributions with the only exception for the F275W filter. For this reason we prefer to apply a more conservative approach and use the peak of the luminosity distribution as a limit for the cluster luminosity function analysis. Only in the case of F275W the 90\\% completeness limit is brighter than the peak magnitude, therefore the latter is adopted as the completeness limit value.\nWe stress that the part of the luminosity distribution leftwards of the peak remains almost identical after the selection cut (check Fig~\\ref{fig:checkmaglimits}) suggesting that the distributions are not affected by our selection criteria and completeness recovery at magnitudes brighter than the peak of the distributions.\n\nWe also tested completeness variations on sub-galactic scales. Our analysis shows that the completeness is worse towards the centre of the galaxy. Outside the inner region (radius larger than 35''), the $V$ band 90\\% completeness value is fainter than the cut at 23.4 mag (see Tab.~\\ref{tab:completeness} and Fig.~\\ref{fig:checkmaglimits}). Similar results are observed in the other filters. Because of the completeness drop at radii smaller than 35'' (1.3 kpc), we have excluded from our analysis this inner region.\n \n\\subsection{SED fitting}\n\\label{sec:sed}\nSources with detection in at least four filters were analyzed via SED fitting algorithms in order to derive physical properties of the clusters. We use uniformly sampled IMF to derive SSP models that include a treatment for nebular continuum and emission lines as described in \\citet{legus2}. \nPutting together the different aperture correction methods, different stellar libraries and different extinction curves, 12 final catalogues are produced with the deterministic fitting method (and will be available online on the LEGUS website\\footnote{\\hyperref[https:\/\/archive.stsci.edu\/prepds\/legus\/]{https:\/\/archive.stsci.edu\/prepds\/legus\/}}). All catalogues use models with solar metallicity for both stars and gas and an average gas covering factor of $50\\%$.\n\nThe analyses and results presented hereafter are obtained using the final catalogue with:\n\\begin{itemize}\n\\item photometry corrected by average aperture corrections;\n\\item Padova evolutionary tracks produced with solar metallicity stellar libraries. \n\\item Milky Way extinction curve \\citep{cardelli1989}.\n\\item \\citet{kroupa2001} stellar initial mass function.\n\\end{itemize}\n\n\\section{Global Properties of the cluster sample}\n\\label{sec:global}\n\\subsection{Final Catalogue}\n\\label{sec:finalcat}\nIn Fig.~\\ref{fig:agemass} we show the ages and masses of the sources classified as class 1, 2 and 3. In the same plot, the completeness limit of 23.4 mag in the $V$ band discussed in the previous section is converted into an estimate of the minimum mass detectable for each age. The line representing this limit follows quite accurately the detected sources with minimum mass at all ages. \nUncertainties on the age and mass values are on average within 0.2 dex. Uncertainties can reach 0.3 dex close to the red supergiant phase (visible as a loop at ages $\\sim10-50$ in Fig.~\\ref{fig:ccd}).\nIn order to study only the cluster population of the grand design spiral and avoid the clusters of NGC 5195, we have neglected the clusters with y coordinate bigger than 11600 (in the pixel coordinates of the LEGUS mosaic) both from Fig.~\\ref{fig:agemass} and from the rest of the analysis. This cut is similar to removing the UVIS pointing centred on NGC 5195. \n\\begin{figure*}\n \\centering\n \\subfigure[]\n {\\includegraphics[width=0.49\\textwidth]{05a_plot_agemass_123.pdf}}\n \\hspace{0mm}\n \\subfigure[]\n {\\includegraphics[width=0.49\\textwidth]{05b_plot_agemass_123_hist.pdf}}\n \\hspace{0mm}\n \\caption{(a) Age-mass diagram for the sources in class 1 and 2 (blue circles) and class 3 (red triangles). The solid black line represents the mass limit as function of ages estimated from the evolutionary tracks assuming a completeness limit of 23.4 mag in V band. The dotted horizontal line shows a mass of 5000 M$_\\odot$. The sources in each age bin were furnished with a small amount of artificial scatter around their respective bin to make the plot easier to read. (b) Histogram of the age distribution of the class 1+2 (blue) and class 3 (red) sources: the total height of each bin gives the total number of sources of the classes 1, 2 and 3. The number of class 3 sources drops fast in the first 10 Myr, such that in the range Log(age) = 7 - 8.5 their number is very small even if the age spanned is more that 100 Myr.}\n\\label{fig:agemass}\n\\end{figure*}\n\nThe different classes are not distributed in the same regions of the plot, with class 3 sources having on average smaller masses and younger ages. Previous studies \\citep{grasha2015,ryon2017,legus2} have shown that our morphological cluster classification is with good approximation also a dynamical classification. \nCompact associations (i.e. class 3) in NGC 628 are on average younger and less massive than class 1 and 2 clusters \\citep{grasha2015}. In addition, the age distributions of class 3 systems suggest that they are more easily and quickly disrupted (see \\citealp{legus2} for details), probably because they are not bound. \\citet{grasha2015,grasha2017}, focusing on the clustering function of clusters in the LEGUS galaxies, have also shown that class 3 clusters behave differently than class 1 and class 2 clusters. These results contribute to the idea that the morphological classification chosen has also some dynamical implication: class 3 sources seem to be short-lived systems (see Fig.~\\ref{fig:agemass}b), possibly already unbound at the time of formation. \n\nFor the remainder of this work we will only analyse class 1 and 2 objects, which we consider to be the candidate stellar clusters, i.e. the gravitationally bound stellar systems that form the cluster population of M51. In total we have 2834 systems classified as class 1 and 2 out of the 10925 sources with 3 sigma detection in at least 4 LEGUS bands. \n\nIn Fig~\\ref{fig:ebv} we show how the recovered extinction changes as a function of cluster age. We see that on average the internal extinction of YSCs changes from $E(B-V)\\sim$0.4 mag at very young ages to 0.2 mag for clusters older than 10 Myr, and is even lower ($E(B-V)\\sim$0.1 mag) for clusters older than 100 Myr. We observe also a large scatter at each age bin, suggesting that the extinction may not only be related to the evolutionary phase of the clusters but also to the region where the cluster is located within the galaxy. \n\\begin{figure}\n \\centering\n \\includegraphics[width=\\columnwidth]{06_ageVSebv.pdf}\n\\caption{Age-extinction diagram for the class 1 and 2 sources in the catalogue (blue circles). The red squares are the median values of E(B-V) in age bins of 0.5 dex. The red lines marks the range between the first and the third quartiles in the same age bins. }\n\\label{fig:ebv}\n\\end{figure}\n\n\n\\subsection{Comparison with previous catalogues}\n\\label{sec:comparison}\nNumerous studies of the cluster population in M51 are available in the literature. We compare our catalogue with recently published ones, testing both our cluster selection and agreement in age and mass estimates. These comparisons will be used when we compare the results of our analyses to values reported in the literature.\nAmong the works that studied the entire galaxy, \\citet{scheepmaker2007}, \\citet{hwang_lee2008} and \\citet{chandar16} used the same $BVI$ data as our work. However, we decided to focus our comparison only on two catalogues for which we have access to estimates of ages and masses.\nThe first catalogue is the one compiled by \\citet{chandar16} (hereafter CH16), using the same $BVI$ observations, plus F658N (H$\\alpha$) filter observations from the same program (GO-10452, PI: S. Beckwith) and WFPC2\/F336W filter ($U$ band) observations (from GO-10501, PI: R.Chandar, GO-5652, PI: R.Kirshner and GO-7375, PI: N.Scoville). The cluster candidates catalogue used in their analysis includes 3812 sources in total (of which 2771 lies in the same area covered by our UVIS observations) and has been made publicly available \\citep{chandar16cat}. \nThe second catalogue is taken from \\citet{bastian2005} (hereafter BA05), covers only the central part of the galaxy and is mainly based on HST observations with the WFPC2 camera. It contains 1150 clusters, 1130 of which are in an area in common with our UVIS pointings.\nThese two catalogues are very different both in terms of coverage and instruments used. For this reason we compare the catalogue produced with LEGUS with each of them separately.\n\n\\subsubsection{Cluster Selection}\nWe first compare the fraction of clusters candidates in common between the catalogs. When doing so, we include class 3 sources in the LEGUS sample, as this class of sources is considered in BA05 and CH16 catalogues. Tab.~\\ref{tab:comparison} collects the results of the comparison. Among the 2771 candidates detected in the same area of the galaxy by CH16 and LEGUS, 1619 (60\\%) systems are in common. We have repeated the comparison in the regions covered by human visual classification in LEGUS, finding a better agreement ($\\sim75\\%$). We take this last value as reference for the common fraction of candidates and justify the drop observed when considering the entire catalogue as given by the ML misplacing class 3 objects into class 4 (as discussed in Section~\\ref{sec:classification}).\nFig.~\\ref{fig:comparison} shows a blow-up of the galaxy with the cluster positions of both catalogues. We selected two different regions, one where the sources have been classified via visual classification and one where only ML is available. We notice that some of the CH16 candidates which do not appear in LEGUS catalogue of classes 1, 2 and 3, have been assigned a class 4. This is true for both regions. Those sources were extracted by the LEGUS analysis but were discarded based on their morphological appearance. The differences between the two catalogues are therefore mostly due to source classification and not by the extraction process.\n\nThe comparison with the BA05 catalogue indicates a poorer agreement, with less than $40\\%$ of sources in common. This discrepancy, observable in Fig.~\\ref{fig:comparison}, may be caused by the difference in the data and in the approach used to extract the clusters. BA05 analysis is based on WFPC2 data, whose resolution is a factor of $\\sim$2 worse than ACS. In addition they do not apply any CI cut, increasing the contamination from stars.\n\n\\begin{table*}\n\\centering\n\\caption{Comparison between our cluster catalogue and the catalogues by \\citet{chandar16} (CH16) and \\citet{bastian2005} (BA05). The number of clusters reported in the columns are: (1) number of cluster candidates detected in the CH16 and B05 catalogues within the same region covered by LEGUS; (2) number of cluster candidates in common between CH16, B05 and LEGUS (class 1,2,3), respectively; (3) number of cluster candidates in the Ch16 and B05 catalogues restricted to the area of the galaxy that has been visually inspected by humans; (4) number of cluster candidates of CH16 and B05 in common with LEGUS class 1, 2 and 3 sources that have been classified by human.\n$^1$ These fractions would increase to 95\\% if we include cluster candidates classified as 4 within the LEGUS catalogue. $^2$ This fraction would increase to 71\\% if class 4 cluster candidates are included. $^3$ The agreement would increase to 67\\% if class 4 objects are considered. \nThe drop in the common fraction of cluster candidates from column (4) to column (2) is likely due to the ML classification misplacing class 3 sources into class 4 (see Section~\\ref{sec:classification}). \nThe low fraction of clusters in common with BA05 is most likely due to a different approach for the cluster extraction analysis and low resolution imaging data. BA05 is based on WFPC2 data (a factor of 2.5 worse spatial resolution with respect to WFC3) and does not include a CI cut. This causes a higher contaminations from unresolved sources, which in our analysis were excluded by the CI selection. \n}\n\\begin{tabular}{lcccc}\n\\hline\nCatalogue\t\t\t&\t\\# clusters\t& \\# clusters\t\t& \\# clusters\t\t& \\# clusters\t\t\\\\\n\\\t\t\t\t& \t(1) \t\t& (2)& (3) & (4)\\\\\n\\hline\n\\hline\nCH16\t\t\t\t\t\t& 2711\t\t\t& 1619 (60\\%)$^1$\t\t\t& 732\t\t\t& 535 (73\\%)$^1$\t\t\\\\\nLEGUS (1,2,3)$_{CH16\\ area}$\t \t& 3240\t\t\t& $-$\t\t\t\t\t& 1294\t\t\t& $-$\t\t\t\\\\\n\\hline\nBA05\t\t\t\t\t\t& 1130\t\t\t& 388 (35\\%)$^2$\t\t\t& 214\t\t\t& 83 (39\\%)$^3$\t\t\\\\\nLEGUS (1,2,3)$_{BA05\\ area}$ \t& 1238\t\t\t& $-$\t\t\t\t\t& 267\t\t\t& $-$\t\t\t\\\\\n\\hline\n\\end{tabular}\n\\label{tab:comparison}\n\\end{table*}\n\n\\begin{figure*}\n \\centering\n\\subfigure{\\includegraphics[width=\\textwidth]{07a_comparison2.png}}\n\\subfigure{\\includegraphics[width=\\textwidth]{07b_comparison3.png}}\n \\caption{Comparison between different cluster catalogues in a spiral arm in the galaxy covered by human visual classification (top) and by machine learning classification (bottom) in LEGUS. Green circles are the clusters candidates in \\citet{chandar16} catalogue. Yellow circles are the cluster candidates in \\citet{bastian2005} catalogue. Red circles are cluster candidates of class 1 and 2 in LEGUS, blue circles are class 3 sources in LEGUS and white dashed circles are sources which have been assigned class 4 in LEGUS. In the bottom right corner an inset shows the position of the zoomed region inside M51.}\n\\label{fig:comparison}\n\\end{figure*}\n\n\\subsubsection{Comparison of Ages and Masses}\n\\begin{figure*}\n \\centering\n \\subfigure{\\includegraphics[width=0.43\\textwidth]{08a_plot_ages_ch16.pdf}}\n \\subfigure{\\includegraphics[width=0.43\\textwidth]{08b_plot_ages_ba05.pdf}}\n \\subfigure{\\includegraphics[width=0.43\\textwidth]{08c_plot_1to1_ch16.pdf}}\n \\subfigure{\\includegraphics[width=0.43\\textwidth]{08d_plot_1to1_ba05.pdf}}\n \\subfigure{\\includegraphics[width=0.43\\textwidth]{08e_plot_masses_ch16.pdf}}\n \\subfigure{\\includegraphics[width=0.43\\textwidth]{08f_plot_masses_ba05.pdf}}\n \\caption{Comparison of ages and masses retrieved from the broadband SED fitting in the LEGUS analysis and in the works of \\citealp{chandar16} (CH16, red, left column) and \\citet{bastian2005} (BA05, green, right column). Only the sources in common between LEGUS and CH16 or BA05 are plotted. First row: ages distributions. Second row: 1--to--1 comparison between ages. The 1--to--1 agreement line is shown in solid blue, the 0.3 dex scatter lines in dashed blue. Third row: masses distributions.}\n\\label{fig:comparison_am}\n\\end{figure*}\nThe comparison of age distributions for the sources in common between LEGUS and CH16 is plotted in Fig.~\\ref{fig:comparison_am} (top left). The age distribution of CH16 has a strong peak for sources younger than log(age\/yr)$=7$ and a subsequent drop in the range $7-7.5$, both of which are not observed in our catalogue. The one--to--one comparison between age estimates in Fig.~\\ref{fig:comparison_am} (middle left) shows a large fraction of clusters with young ages in CH16 which have a wide age spread in the LEGUS catalogue. More in general, the differences in the age estimates are mostly caused by the different broadband combinations used in fitting the data, as already noticed in CH16. In addition to the ``standard'' $UBVI$ filter set used for SED fitting of both CH16 and our LEGUS catalogues, we use an extra $UV$ broadband while CH16 use the flux of the narrow band filter centred on the H$\\alpha$ emission line, from an aperture of the same size as the broadband ones. Both approaches aim at breaking the age-extinction degeneracy weighting different information. The LEGUS standard approach is to use two datapoints below the Balmer break ($\\lambda < 4000$) which give a stronger constraint on the slope of the spectrum, and thus, extinction. The approach used by CH16 is to use the detection of H$\\alpha$ emission from gas ionised by massive stars to determine the presence of a very young stellar population in the cluster. From Fig.~\\ref{fig:comparison_am} (middle left) we observe that the two methods agree within 0.3 dex in $\\sim50$\\% of the cases. \nThe correlation between the ages derived in the two methods is confirmed by a Spearman's rank correlation coefficient $r_s=0.7$ with a p-value: $10^{-234}$.\nWe will address in a future work (Chandar et al. in prep.) the systematics and differences in the two methods. In this work we will take into account the differences observed in the age distributions when discussing and comparing our results to those available in the literature.\n\nThe mass distributions (Fig.~\\ref{fig:comparison_am}, bottom left) show a more similar behaviour, with a broad distribution and a decrease at low masses caused by incompleteness. Note that CH16 retrieve higher mass values at the high-mass end of the distribution. This difference can be important in the study of the mass function shape (Section~\\ref{sec:massfunc}).\n\nThe comparison of age and mass distributions for the sources in common between LEGUS and BA05 is shown in the right column of Fig.~\\ref{fig:comparison_am}.\nSince the youngest age assigned by BA05 is log(age\/yr)=6.6, for the sake of the comparison in Fig.~\\ref{fig:comparison} (top right) we have assigned log(age\/yr)=6.6 age to all the sources that in our catalogue are younger. The general trends of the distributions looks similar, but the one--to--one comparison in Fig.~\\ref{fig:comparison} (middle right) reveals that the two methods agree within 0.3 dex in $\\sim50\\%$ of the cases. \nThe correlation found with a Spearman's rank test is $r_s=0.5$. A p-value of $10^{-24}$ confirms that this correlation is not random, but the moderate value of $r_s$ is caused by the difference in the age distribution observed in Fig.~\\ref{fig:comparison} (middle right). \nBA05 use very different data from our own and allow fits with $BVI$ bands only, with large uncertainties on the recovered properties. For example, the large cloud of systems that sit in the upper left part of the plot has been assigned younger ages in our catalogue. Also in this case we can conclude that the differences in age determinations are mostly caused by the different fitting approach, with our catalogue having more information to break the age-extinction degeneracy. The mass distributions (Fig.~\\ref{fig:comparison_am}, bottom right) show the same overall shape, with the BA05 distribution shifted by 0.2 dex to higher values of masses.\n\nIn general, for both catalogues, we notice that differences in the derived properties can be also caused by differences in the stellar templates adopted, which are different for all catalogues (CH16 uses \\citealp{bruzual_charlot2003} models, while BA05 uses updated GALEV simple stellar population models from \\citealp{schulz2002} and \\citealp{anders2003}). We will use the differences outlined among these previously published catalogues and ours when we will discuss the results of our analyses.\n\n\n\\subsection{Cluster Position as a function of age}\n\\label{sec:position}\nIn order to understand where clusters form and how they move in the dynamically active spiral arm system of M51 we plot the position of the clusters inside the galaxy in Fig.~\\ref{fig:positions}. The sample is divided in age bins ($1-10$, $10-100$, $100-200$ and $200-500$ Myr). \nClusters in our sample are mostly concentrated along the spiral arms. This trend is particularly clear for the very young clusters (age $<10$ Myr) but can also be spotted in the ranges $10-100$ and $100-200$ Myr. In general we observed that young sources are clustered. Moving to older sources the spatial distribution becomes more spread, but it can still be recognized that sources are more concentrated along the spiral arms. In the last age bin, probing clusters older than 200 Myr, the number of available sources is much smaller and is therefore hard to define a distribution, although the sources appear to be evenly spread across the area covered by observations. The strength and age-dependency of the clustering will be further investigated in a future paper (Grasha et al, in prep).\n\\begin{figure*}\n \\centering\n \\subfigure\n {\\includegraphics[width=0.4\\textwidth]{09a_positions1.pdf}}\n \\subfigure\n {\\includegraphics[width=0.4\\textwidth]{09b_positions2.pdf}}\n \\subfigure\n {\\includegraphics[width=0.4\\textwidth]{09c_positions3.pdf}}\n \\subfigure\n {\\includegraphics[width=0.4\\textwidth]{09d_positions4.pdf}}\n \\caption{Position of the clusters divided in age bins in the ranges 1-10 Myr (left), 10-100 Myr (middle-left), 100-200 Myr (middle-right) and 200-500 Myr (right). The central circle of radius 1.3 kpc encloses the region where detection is poorer (see text). The UVIS footprint, restricted to y values below 11600, is overplotted as a grey solid line.}\n\\label{fig:positions}\n\\end{figure*}\n\nThe lack of age gradient as a function of distance from the spiral arm observed in Fig. 8 is in agreement with the detailed study of azimuthal distances of clusters as a function of their ages collected in a forthcoming paper (Shabani et al., in prep.) where the origin of spiral arm and dynamical evolution is investigated.\nThe observed trend has been predicted by \\citet{dobbs2010} which, modeling a spiral structure induced by tidal interactions find that clusters of different ages tend to be found in the same spiral arm without a defined age gradient. \nIn a more recent numerical work, \\citet{dobbs2017} analyse the evolution of stellar particles in clustered regions, i.e. simulated star clusters within spiral fields. They observe that up to the age range they are able to follow (e.g. 200 Myr) their simulated clusters are mainly distributed along the spiral arms. The trend observed in M51 is thus compatible with that found in \\citet{dobbs2017} simulations. Simulations on the evolution of M51 \\citep[e.g][]{dobbs2010_m51} suggest that the interaction with the companion galaxy, started $\\sim300$ Myr ago, is responsible for creating or strengthening the spiral arms and may have helped keep the old clusters we see now fairly close to the arms.\n\nFrom Fig.~\\ref{fig:positions} we clearly see that our detection is very poor in the centre of the galaxy where the bright background light of the diffuse stellar population is much stronger than in the rest of the galaxy. This effect could explain why we do not detect sources older than 10 Myr (i.e. when cluster light starts to fade), causing a drop in the completeness limit, as already pointed out in Section~\\ref{sec:completeness}. For this reason we ignore the clusters within 35'' (1.3 kpc) from the centre of the galaxy from the following analyses. \n\n\n\n\n\\subsection{Luminosity function}\n\\label{sec:lumfunc}\n\nThe luminosity function is intrinsically related to the mass function (luminosity is proportional to mass, with a dependence also on the age) but it is an observed quantity, and therefore, like the colour-colour diagrams, available without any assumption of stellar models and without any SED fitting. \nThe luminosity function of YSCs is usually described by a power law function $dN\/dL\\propto L^{-\\alpha}$, with an almost universal index close to $\\alpha\\approx -2$ as observed in local spiral galaxies (e.g. \\citealp{larsen2002,degrijs2003}, see also the reviews by \\citealp{whitmore2003} and \\citealp{larsen2006b}).\n\nWe analyse the cluster luminosity function by building a binned distribution with the same number of objects per bin, as described in \\citet{maiz2005} and performing a least-$\\chi^2$ fitting.\nThe errors on the data are statistical errors given by $\\sigma_{\\textrm{bin}}=\\sqrt{\\frac{n_{\\textrm{bin}}(n_\\textrm{tot}-n_\\textrm{bin})}{n_\\textrm{tot}}}$, where $n_\\textrm{bin}$ is the number of sources in each bin and $n_\\textrm{tot}$ is the total number of sources. \nThe results of the fits are listed in Tab.~\\ref{tab:lumfit} and plotted in Fig.~\\ref{fig:lumfit}. \nWe have fitted the data up to the completeness limits described in Section~\\ref{sec:completeness}.\nThe function is fitted with both a single and a double power law (PL). The single PL fit gives slopes close to a value of $\\alpha=-2$,\nhowever, for all filters, the double power law results in a better fit, as the $\\chi^2_{red}$ in this second case is always lower. \n\\begin{table*}\n\\centering\n\\caption{Results of the fit of the binned luminosity functions.}\n\\begin{tabular}{|c|c|c|c|c|c|c|c|c}\n\\hline\n\\multicolumn{1}{|c|}{Filter} \t\t& \\multicolumn{1}{c|}{$\\textrm{Mag}_{\\textrm{cut}}$} \t& \\multicolumn{2}{l|}{Single PL fit} \t& \\multicolumn{4}{l|}{Double PL fit} & \\multicolumn{1}{l|}{Cumulative fit}\\\\\n\\multicolumn{1}{|c|}{} \t\t& \\multicolumn{1}{c|}{} \t\t& \\multicolumn{1}{c|}{$\\alpha$} & \\multicolumn{1}{c|}{$\\chi^2_{red.}$}\t\t& \\multicolumn{1}{c|}{$\\alpha_1$}\t& \\multicolumn{1}{c|}{$\\textrm{Mag}_{\\textrm{break}}$} \t& \\multicolumn{1}{c|}{$\\alpha_2$} & \\multicolumn{1}{c|}{$\\chi^2_{red.}$} & \\multicolumn{1}{c|}{$\\alpha$}\t\\\\\n\\hline\n\\hline\nF275W &\t21.71\t& $1.84\\ _{\\pm 0.04}$\t& 2.17 &\t$1.57\\ _{\\pm 0.06}$ &\t$19.39\\ _{\\pm 0.25}$ & \t$2.24\\ _{\\pm 0.10}$ & 0.85 & $2.08\\ _{\\pm 0.01}$ \\\\\nF336W &\t22.00\t& $1.89\\ _{\\pm 0.04}$\t& 1.80 &\t$1.67\\ _{\\pm 0.05}$ &\t$19.49\\ _{\\pm 0.26}$ & \t$2.32\\ _{\\pm 0.12}$ & 0.92 & $2.10\\ _{\\pm 0.01}$ \\\\\nF435W &\t23.25\t& $1.99\\ _{\\pm 0.03}$\t& 1.72 &\t$1.74\\ _{\\pm 0.04}$ &\t$20.97\\ _{\\pm 0.17}$ & \t$2.41\\ _{\\pm 0.09}$ & 0.80 & $2.17\\ _{\\pm 0.01}$ \\\\\nF555W &\t23.25\t& $2.02\\ _{\\pm 0.03}$\t& 1.70 &\t$1.79\\ _{\\pm 0.03}$ &\t$20.83\\ _{\\pm 0.17}$ & \t$2.48\\ _{\\pm 0.10}$ & 0.86 & $2.18\\ _{\\pm 0.01}$ \\\\\nF814W &\t22.25\t& $2.04\\ _{\\pm 0.04}$\t& 2.31 &\t$1.60\\ _{\\pm 0.07}$ &\t$20.73\\ _{\\pm 0.13}$ & \t$2.40\\ _{\\pm 0.08}$ & 1.03 & $2.28\\ _{\\pm 0.01}$ \\\\\n\\hline\n\\end{tabular}\n\\label{tab:lumfit}\n\\end{table*}\n\\begin{figure*}\n \\centering\n \\subfigure[binned]\n {\\includegraphics[width=0.48\\textwidth]{10a_plot_lumfunc_DP_new.pdf}}\n \\hspace{0mm}\n \\subfigure[cumulative]\n {\\includegraphics[width=0.48\\textwidth]{10b_plot_lumcumfit_DP_new.pdf}}\n \\caption{Binned (a) and cumulative (b) luminosity functions. Fit results are reported in Tab. \\ref{tab:lumfit}. The curves are for different filters, from UV band at the top to I band at the bottom. Errorbars in (a) are of the same size of the markers. In both panels the vertical lines mark the completeness limit in each filter.}\n\\label{fig:lumfit}\n\\end{figure*}\n\nSimilar results were found by \\citet{haas2008} using a cluster catalogue based only on $BVI$ photometry.\nThey found that the low-luminosity part of the function could be fitted by a shallow power-law, with slopes in the range $\\sim1.7-1.9$, while the high-luminosity end was steeper, with slopes $\\sim2.3-2.6$. We similarly found that the low luminosity part of the function is shallower ($\\alpha\\sim1.6-1.8$) than the high luminosity part ($\\alpha\\sim2.4-2.5$). \nIn both analyses a double power law is a better fit of the luminosity function in all filters. As suggested by \\citet{gieles2006}, a broken power law luminosity function suggests that also the underlying mass function has a break.\nThe possibility that the underlying mass function is truncated is further explored with the study of the mass function in Section~\\ref{sec:massfunc} and via Monte Carlo simulations in Section~\\ref{sec:montecarlo}.\n\nWe compared the binning fitting method with the one presented in \\citet{bastian2012}, involving the use of cumulative functions. In the case of a single power law behaviour, the two functions are expected to show the same shape. The cumulative function is given by $y_{\\textrm{c}}(m)=\\left(1 - \\frac{k}{n_{\\textrm{data}}}\\right)$, where $k$ is the index of the object of magnitude m in the sorted array containing the magnitudes and $n_{\\textrm{data}}$ is the length of the array. In case of a simple power law, it has a slope $\\alpha_c=\\alpha_b-1$ which can be directly compared to the slope $\\alpha_b$ of the binned function. Also in this case a least $\\chi^2$ fit is performed. In the cumulative distribution no error is associated with the data, therefore the fit is made with a linear function in the logarithmic space, assigning the same uncertainty to all points. The errors on the fitted parameters have been estimated via a bootstrapping technique: 1000 Monte Carlo realisations of the distribution are simulated, where the luminosity of each cluster is changed using uncertainties normally distributed around $\\sigma_{mag}=0.35$ mag. Each realisation was then fitted in the same way as the original one. The standard deviation of the 1000 values recovered for each parameter gave the final uncertainty associated with the recovered slopes.\nResults and slopes are collected in Fig.~\\ref{fig:lumfit}(b) and Tab.~\\ref{tab:lumfit}. We converted $\\alpha_c$ into $\\alpha_b$ in Tab.~\\ref{tab:lumfit}, for an easier comparison with the binned function.\nThe trends observed in the analysis of the binned distributions are traceable in the cumulative function as well. In particular, we observe that the single PL fit is not a good description of the bright end of the cumulative distributions in all the filters (Fig.~\\ref{fig:lumfit}b). \nAlso with the cumulative function, the brightest sources fall below the expected curve of a single power law distribution, sign of a break in the luminosity function and therefore also in the underlying mass function.\nWe note that the fit of the cumulative functions results in steeper slopes than the ones recovered with the binned distributions. This discrepancy caused by the differences in the two techniques is discussed at length in \\citet{legus2}.\n \n\n\n\n\\subsection{Mass Function}\n\\label{sec:massfunc}\nThe results obtained with the analysis of the luminosity function can be further explored with the study of the properties derived from the SED fit, i.e. the mass and the age distributions.\nIn the following analyses we use a mass-complete sample, by selecting only clusters above $5000\\ \\textrm{M}_{\\odot}$. This value has been chosen in order to avoid low mass sources, affected by inaccuracies in the SED fitting and by stochastical sampling of the stellar IMF (see e.g. the comparison between deterministic and Bayesian fitting of cluster SEDs in Figure 15 of \\citealp{slug}).\nThe age-mass plot of Fig. \\ref{fig:agemass} suggests that we are complete in recovering sources more massive than $5000\\ \\textrm{M}_{\\odot}$ only up to 200 Myr. At older ages, even sources more massive than 5000 M$_\\odot$ can fall below our magnitude detection limit. \nOur mass-limited complete sample therefore contains sources with M $>5000\\ \\textrm{M}_{\\odot}$ and ages $< 200$ Myr.\n\nThe cluster mass function is expected to evolve from a cluster initial mass function (CIMF), usually assumed as a power law $dN\/dM \\propto M^{\\beta}$ with a $\\beta=-2$ slope. This slope is interpreted as the sign of the formation of clusters from a turbulent hierarchical medium \\citep{elmegreen2010b}. The initial function is then expected to evolve due to cluster evolution and disruption. \n\nThe cluster mass function of our sample is plotted in Fig.~\\ref{fig:massfit}, where bins of equal number of sources were used.\nWe recover a shape which is well fitted with a single power law of slope $-2.01 \\pm0.04$ ($\\chi^2_{red.}$ of 1.6), even if a double power law with a steeper high mass slope fits better the function ($\\chi^2_{red.}$ of 1.1, see Tab.~\\ref{tab:massfit}). \n\\citet{gieles2009} and CH16 found similar slopes, $\\beta=-2.09 \\pm0.09$ and $\\beta=-1.97 \\pm0.09$ respectively, considering only clusters in the age range from 10 to 100 Myr. Restricting to the same age range, we find a consistent value of $\\beta=-2.03 \\pm0.04$ ($\\chi^2_{red} = 0.77$, see Fig.~\\ref{fig:massfunc_ages}).\n\nAs done in the analysis of the cluster luminosity function, we also plot the mass function in a cumulative form (Fig.~\\ref{fig:massfit}) and fit it with a pure power law. As already observed for the luminosity functions, the cumulative mass distributions show a steepening at the high-mass end. As observed in \\citet{bastian2012} and \\citet{legus2}, while the equal number of object binning technique is statistically more robust, it is insensitive to small scales variations, like the dearth of very massive clusters. The cumulative form is therefore more appropriate to study the high-mass end of the mass (and luminosity) function.\n\\begin{table}\n\\centering\n\\caption{Values derived by the fit of the mass function with a least-$\\chi^2$ fitting of the binned function. Fits have been performed considering a low mass cut of either 5000 or $10^4$ M$_\\odot$, as indicated in the second column.}\n\\begin{tabular}{lccccc}\n\\hline\nMethod\t& \tM cut \t& $-\\beta_1$\t\t\t& $M_{break}$\t& $-\\beta_2$ \t& $\\chi^2_{red.}$\\\\\n\\\t\t&\t[M$_\\odot$]\t& \\\t\t\t\t\t& [M$_\\odot$]\t\t& \\\t\t\t& \\\t\t\t\\\\\n\\hline\n\\hline\nSingle PL\t&\t5000\t\t& $2.01_{\\pm0.04}$\t\t& $-$\t\t\t& $-$\t\t\t& 1.6\\\\\nDouble PL\t&\t5000\t\t& $1.52_{\\pm0.12}$\t\t& $1.5\\times10^4$\t& $2.31_{\\pm0.09}$\t& 1.1\\\\\n\\hline\nSingle PL\t&\t$10^4$\t& $2.19_{\\pm0.07}$\t\t& $-$\t\t\t& $-$\t\t\t& 1.6\\\\\nDouble PL\t&\t$10^4$\t& $1.71_{\\pm0.39}$\t\t& $1.8\\times10^4$\t& $2.36_{\\pm0.13}$\t& 1.5\\\\\n\\hline\n\\end{tabular}\n\\label{tab:massfit}\n\\end{table}\n\\begin{table}\n\\centering\n\\caption{Values derived by the fit of the cumulative mass function with a maximum-likelihood fit. For a description of the values $N_0$ and $M_0$ see Eq.~\\ref{eq:cumulative}. In the last row the result of the fit of the GMC population is reported.}\n\\begin{tabular}{lccccc}\n\\hline\nMethod\t\t& M cut \t& Age \t&$-\\beta$\t\t\t\t&$N_0$\t\t\t& $M_0$ \\\\\n\\\t\t\t& [M$_\\odot$]\t& [Myr]\t& \\\t\t\t\t\t& \\\t\t\t\t& [$10^5$ M$_\\odot$] \\\\\n\\hline\n\\hline\nSingle PL\t\t& 5000\t& 1-200\t& $2.30_{\\pm0.03}$\t\t& $-$ \t\t\t& $-$ \\\\ \nTruncated \t& 5000\t& 1-200\t& $2.01_{\\pm0.02}$\t\t& $66_{\\pm6}$\t\t& $1.00_{\\pm0.12}$\\\\\nTruncated\t\t& 5000\t& 1-10\t& $2.12_{\\pm0.22}$\t\t& $10_{\\pm7}$\t\t& $0.56_{\\pm0.08}$\\\\\nTruncated\t\t& 5000\t& 10-100\t& $1.97_{\\pm0.06}$\t\t& $43_{\\pm15}$\t& $0.91_{\\pm0.16}$\\\\\nTruncated\t\t& 5000\t& 100-200\t& $2.01_{\\pm0.05}$\t\t& $28_{\\pm4}$\t\t& $1.15_{\\pm0.27}$\\\\\n\\hline\nSingle PL\t\t& $10^4$\t& 1-200\t\t& $2.67_{\\pm0.03}$\t\t& \\ \t\t\t\t& \\ \\\\ \nTruncated\t\t& $10^4$\t& 1-200\t\t& $2.34_{\\pm0.03}$\t\t& $22_{\\pm10}$\t& $1.34_{\\pm0.24}$\\\\\n\\hline\nGMC pop \t\t& $-$\t& $-$\t& $2.36_{\\pm0.16}$ \t& $12_{\\pm5}$\t\t& $160_{\\pm32}$ \\\\\n\\hline\n\\end{tabular}\n\\label{tab:massfit2}\n\\end{table}\n\nIn order to test the hypothesis of a mass truncation, we have fitted the cumulative distribution with the \\texttt{IDL} code \\texttt{mspecfit.pro}, implementing the maximum-likelihood fitting technique described in \\citet{rosolowsky2007}, commonly used for studying the mass functions of GMCs \\citep[e.g.][]{colombo2014a}. The code implements the possibility of having a truncated power law mass function, i.e.\n\\begin{equation}\n\\label{eq:cumulative}\nN(M'>M)=N_0\\left[\\left(\\frac{M}{M_0}\\right)^{\\beta+1}-1\\right],\n\\end{equation}\nwhere $M_0$ is the maximum mass in the distribution and $N_0$ is the number of sources more massive than $2^{1\/(\\beta+1)}M_0$, the point where the distribution shows a significant deviation from a power law (for the formalism, see \\citealp{rosolowsky2005}).\nA value of $N_0$ bigger than $\\sim1$ would indicate that a truncated PL is preferred over a simple one. On the other hand, $N_0<1$ would mean that the truncation mass is not constrained and that a single power law is a good description of the distribution. The resulting parameters of the fit for our sample, considering normally distributed 0.1 dex errors on the masses, are collected in Tab.~\\ref{tab:massfit2}. \nThe resulting $N_0 = 66 \\pm6$ suggests that the fit with a truncated function,with $M_0=10^5$ M$_\\odot$, is preferred over the simple PL. The best fit for the slope is $\\beta=-2.01 \\pm0.02$.\n\\begin{figure}\n\\centering\n\\includegraphics[width=\\columnwidth]{11_plots_massfit_cat_bestm_all.pdf}\n\\caption{Binned (top) and cumulative (bottom) mass function. The solid lines are the best fits with a single power law, the dashed line in the case of the cumulative function shows the slope $-2$ for comparison. The binned function is steeper because the slope of the cumulative function in a simple power law case is $B=\\beta+1$, where $\\beta$ is the slope of the original function (see Eq.~\\ref{eq:cumulative}).}\n\\label{fig:massfit}\n\\end{figure}\n\nIn order to test for possible incompleteness at masses close to 5000 M$_\\odot$, we repeated the analysis of the mass function using a mass cut of $10^4$ M$_\\odot$.\nResults are collected in Tab.~\\ref{tab:massfit} and \\ref{tab:massfit2}. Different lower limits at the low mass yield steeper than $-2$ power laws but consistent truncation masses. \nThe binned function is well fitted with a single power law with $\\beta=-2.19 \\pm0.07$ ($\\chi^2_{red.}=1.55$). The maximum likelihood fit of the cumulative function gives $\\beta=-2.34 \\pm0.03$, $M_0= (1.34 \\pm 0.24) \\times10^5$ M$_\\odot$\\ and $N_0=22 \\pm 10$, thus a truncation is still statistically significant. \n\nIt has been reported in the literature that the YSC mass function is probably better described by a Schechter function with a $\\beta=-2$ slope and a truncation $M_c$ at the high-mass end. In the case of M51, \\citet{gieles2006} found that a Schechter function with $M_c = 10^5$ M$_\\odot$\\ would reproduce closely the luminosity function observed. With a very different approach, \\citet{gieles2009} derived $M_c = (1.86 \\pm0.52) \\times 10^5$ M$_\\odot$\\ from the analysis of an evolving mass function. Both results are consistent with our results. \n\n\\citet{chandar2011,chandar16} find that a simple PL is a good description of the YSC mass function in M51, however they only considered a binned MF. Different mass estimates for high-mass clusters, as noted in Fig.~\\ref{fig:comparison_am}, could produce differences in the mass function slopes. Nevertheless, we retrieve the same results of CH16 if a binned function is used. We notice that the bin containing the most massive clusters encompasses the whole range in masses where the truncation mass is found (see Fig.~\\ref{fig:massfit}). Thus binning techniques that use equal number of objects are therefore unable to put a constraint on the truncation. \n\nThe truncation mass we recover is smaller but similar to what was found in other spirals, like M 83 ($M_c=(1.60 \\pm0.30) \\times10^5$ M$_\\odot$, \\citealp{adamo2015}), NGC 1566 ($M_c=2.5 \\times10^5$ M$_\\odot$, \\citealp{hollyhead2016}) and NGC 628 ($M_c=(2.03 \\pm0.81) \\times 10^5$ M$_\\odot$, \\citealp{legus2}). On the other hand, some galaxies still exhibit a truncated mass function but with very different truncation masses. In M31, for example, \\citet{johnson2017} found a remarkably small truncation mass of $\\sim10^4$ M$_\\odot$, while the Antennae have a MF that exhibits a PL shape that extends up to masses larger than $10^6$ M$_\\odot$\\ \\citep{whitmore2010}. These differences spanning orders of magnitudes suggest that the maximum cluster mass in galaxies must be determined by the internal (gas) properties of the galaxies themselves. \\citet{johnson2017} suggested that $M_c$ should scale with the $\\Sigma_{\\textrm{SFR}}$.\nDifferences in the recovered truncation mass have also been found within the same galaxy (e.g. \\citealp{adamo2015}). We will investigate possible environmental dependencies of the mass function properties of M51 in a follow up work (Paper II).\n\n\n\n\\subsubsection{Comparison With GMC Masses}\n\\label{sec:GMC}\nWe compare our cluster mass function with the mass function of the GMCs in M51 from the catalogue compiled and studied in \\citet{colombo2014a}. Clusters are expected to form out of GMCs, via gravitational collapse and fragmentation, and therefore the mass distribution of the latter can in principle leave an imprint on the mass distribution of YMCs.\n\nThe mass function of GMCs in M51 steepens continuously going from low to high masses, and cannot be described by a single power law \\citep{colombo2014a}, as is instead the case for other galaxies like LMC, M33, M31 and the Milky Way \\citep{wong2011,gratier2012,rosolowsky2005}. We perform a fit of GMC masses with the same code described in the previous section, up to a lower limiting mass of $10^6$ M$_\\odot$\\ (discussed in Section 7.2 and 7.3 on the mass functions in \\citealp{colombo2014a}). The resulting best value for the slope and the maximum mass are $\\beta=-2.36 \\pm0.16$ and $M_0=(1.6 \\pm0.3) \\times10^7$ M$_\\odot$. \nThe value of $M_0$ implies a truncation mass which is $\\sim100$ times bigger in the case of GMCs similar to what has been observed in M83 by \\citet{freeman2017}. \nThe mass function of GMCs looks steeper than the one of the clusters, within a 3$\\sigma$ difference. Analysing the mass function of simulated GMCs and clusters, \\citet{dobbs2017} found the opposite trend of a steeper function in the case of clusters. Part of the difference between the simulated and observed trends can be due to the different regions covered within the two surveys: the PAWS survey from which the GMC data are derived covers only the central part of the galaxy, while our clusters also occupy more distant regions from the centre.\nThe study of the mass function in different regions of M51 in Paper II will enable us to compare the CMF with the GMC one on local scales, testing closely the link between GMC and cluster properties.\n \n\\subsubsection{Evolution of the Mass Function}\n\\label{sec:massage}\nCluster disruption affects the mass function and could, in principle, modify its shape: for this reason we study the function in different age bins. \nIn order to be able to see how significant the disruption is, we look at the evolution of the CMF normalized by the age range (i.e. $dN\/dMdt$).\nIn case of constant star formation and no disruption the mass functions should overlap. \nCluster disruption can in principle affect the mass function in different ways according to the disruption model considered. \n\nTwo main empirical disruption scenarios have been proposed in the literature and they differ in the dependence with the cluster mass. A first model, firstly developed to explain the age distribution of clusters in the Antennae galaxies (see \\citealp{fall2005} and \\citealp{whitmore2007}), proposes that all clusters lose the same fraction of their mass in a given time. This implies that the disruption time of clusters is independent on the cluster mass and we therefore call this model mass-independent disruption (MID). It is characterized by a power-law mass decline and therefore by a disruption rate which depends linearly on the mass \\citep{fall2009}, i.e.\n\\begin{equation}\nM(t)\\propto t^\\lambda, \\ \\ \\ \\ \\frac{dM}{dt}\\propto M\n\\end{equation}\nOn the other hand, the mass dependent disruption (MDD) time scenario assumes that the lifetime of a cluster depends on its initial mass, with a relation $t_{dis}\\propto M^k$ (with $k=0.65$, i.e. less massive clusters have shorter lifetimes).\nInitially suggested by \\citet{boutloukos2003} considering only instantaneous disruption to explain the properties of the cluster populations in the SMC, M33 and M51, this model has been updated to account also for gradual mass-loss in \\citet{lamers2005}. This model is characterized by a disruption rate which depends sub-linearly on the mass as:\n\\begin{equation}\n\\frac{dM}{dt}\\propto M^{1-k}\n\\end{equation}\nThe two scenarios predict different evolutions for the cluster mass function \\citep[e.g.][]{fall2009}. In the MID model, the mass function shape is constant in time, it only shifts to lower masses due to all clusters losing the same fraction of mass. In the MDD model, instead, low-mass clusters have shorter lifetimes and this results in a dearth of clusters at the low-mass end of the function, as the time evolves.\n\nIn Figure~\\ref{fig:massfunc_ages} we observe that the normalised CMF at ages below 10 Myr is detached from the CMFs of the other two age bins, suggesting a stronger drop in the number of sources. In the age range $10-100$ Myr, compared to the range $100-200$ Myr, the main difference between the two CMF is seen at low masses as a bend in the CMF of the oldest clusters, i.e. $100 - 200$ Myr. This trend would suggest a shorter disruption time for low mass clusters, however as can be seen in Fig~\\ref{fig:agemass}, at these ages also incompleteness could start affecting the data. So the flattening could be the result of both mass dependent disruption time and incompleteness. On the other hand, at high masses the functions seem to follow each other quite well. \n\\begin{figure}\n\\centering\n\\includegraphics[width=1.\\columnwidth]{12_plots_massfit_cat_bestm_all_ages_new.pdf}\n\\caption{Mass function divided in age bins and normalized by the age range in each bin. The black dotted line is the low-mass limit of 5000 M$_\\odot$. The shift in normalisation between the young function (blue) and the others suggest that cluster disruption is already happening between 10 and 100 Myr. The old function (red) flattens at low masses, but it is difficult to separate the effects of disruption and incompleteness.}\n\\label{fig:massfunc_ages}\n\\end{figure}\n\nEach function is fitted with a least-$\\chi^2$ approach. Single power laws are fitted and the resulting slopes for the age bins $1-10$, $10-100$ and $100-200$ Myr are $\\beta=-2.02 \\pm0.11$, $-2.03 \\pm0.04$ and $-1.85 \\pm0.06$, respectively. We can compare these values with the results of CH16 and \\citet{gieles2009} as both of them studied the mass function in age bins. CH16 found a slope $\\beta=-2.06 \\pm 0.05$ for sources younger than 10 Myr and $\\beta=-1.97 \\pm 0.09$ for sources in the range 10-100 Myr. For older sources they consider a bin with ages in the range 100-400 Myr finding a slope of $\\beta=-2.19 \\pm 0.06$. \nUsing age bins of $4-10$ Myr, $10-100$ Myr and $100-600$ Myr, \\citet{gieles2009} found slopes of $\\beta=-2.08 \\pm 0.08$, $-2.09 \\pm 0.09$ and $-2.76 \\pm 0.28$ respectively, using the cluster catalogue of B05. While up to 100 Myr those values are comparable with what we find, at old ages their results seem to strongly deviate from our own. A first reason for this deviation may be the smaller size of our last age bin, which extends only to 200 Myr and therefore neglects older sources. \nHowever, a more likely explanation can be found in the definition of the minimum mass considered in each age range: because of our completeness limit, we always consider sources more massive than 5000 M$_\\odot$, while the cut of the older bin in \\citet{gieles2009} is 6$\\times10^4$ M$_\\odot$\\ and in CH16 is $\\approx 10^4$ M$_\\odot$. In both cases the low-mass part of the function is not considered in the fit, thus their fit may be more sensitive to the presence of a bend in the form of a truncation. As seen in Tab.~\\ref{tab:massfit}, fitting only the high mass part of the CMF results in stepper slopes also in our catalogue, even if a shorter age range is used.\n\nFitting cumulative instead of binned distributions with the maximum-likelihood fit with the \\texttt{mspecfit.pro} code yields slopes and truncation masses collected in Tab.~\\ref{tab:massfit2}. Results for age bins $10-100$ Myr and $100-200$ Myr are very similar to the results found for the whole population. In both age ranges the presence of a truncation ($N_0 >>1$) is statistically significant. The CMF in the age range $1-10$ Myr has a fitted $M_0$ which is a factor 2 smaller. The statistical significancy $N_0$ of the latter fit is, within uncertainties, not much larger than 1 ($N_0=10 \\pm7$). This result seems driven by size--of--sample effects. Uncertainties in this last case are larger because the sample is small, counting only 140 clusters, compared to the other two age bins hosting more than 500 clusters each. These uncertainties prevent to statistically test the truncation for the mass function in the bin $1-10$ Myr.\n\n\n\n\\subsection{Age Function}\n\\label{sec:agefunc}\nWe can investigate the cluster evolution analysing the age distributions of the clusters. The YSC age function is determined by the star (and cluster) formation history (SFH and CFH) convolved with cluster disruption.\n\nIn first approximation, the SFH of spiral galaxies can be assumed constant for extended periods, unless external perturbations (like interactions, minor, or major mergers) change the condition of the gas in the galaxy. \nYSC disruption is usually inferred by changes in the number of clusters as a function of time, assuming that the SFH has been constant. In the presence of enhancement in SF, the change brought by the increasing SFR can be misinterpreted as disruption. Thus it is of fundamental importance to know the recent SFH of the host galaxy. \nThe easiest assumption of a constant star formation rate allows a very straightforward interpretation of the age function which it is not necessarily true. In the case of M51, we know it is an interacting system and that tidal interactions can enhance the star formation \\citep{pettitt2017}. Many simulations of the M51 evolution have suggested a double close passage of the companion galaxy, the older approximately $400-500$ Myr ago and a more recent one $50-100$ Myr ago \\citep[see][]{salo2000,dobbs2010_m51}. Whether the enhancement of star formation during the close passages with the interacting galaxy has a visible impact on the age function is difficult to assess, without an accurate star formation history. Our analysis is limited to an age range $<200$ Myr, hence we expect our analysis to be only partially affected.\n\n\n\\begin{figure}\n\\centering\n\\includegraphics[width=\\columnwidth]{13_plots_agefunc.pdf}\n\\caption{Age function of the cluster catalogue, comparing the effect of different bin widths, namely widths of 0.5 dex (blue squares), 0.6 dex (orange triangles) and 0.7 dex (red circles). The grey-shaded area marks the ages at which incompleteness causes a steepening of the slope, preventing the study of the function.}\n\\label{fig:agefit}\n\\end{figure}\nWe build an age function dividing the sample in age bins of 0.5 dex width and taking the number of sources in each bin. This number is normalised by the age range spanned by each bin (Fig.~\\ref{fig:agefit}). Points have been fitted with a simple power law $dN\/dt \\propto t^{-\\gamma}$ up to $\\textrm{log(age\/yr)}=8.5$. After that the incompleteness strongly affects the shape of the function, which starts declining steeply. The resulting slope is $\\gamma=0.38 \\pm0.06$, smaller than reported by CH16 who found $\\gamma$ in the range $0.6-0.7$ using a set of mass selections and age intervals. \nNote that, in their cluster selection, CH16 do not remove clusters in the internal part of the galaxy, which happen to have a much steeper slope. However the main difference can probably to be attributed to the different age modeling of the two catalogues. As can be seen from Fig. \\ref{fig:comparison} the ages of the CH16 catalogue have a prominent peak between logarithmic ages of 6.5 and 7. The resulting cluster ages are therefore younger on average and result in a steeper decline.\n\nA source of uncertainty in the recovered slope of the age function is related to the binning of the data. The distribution of the ages is discrete, and therefore binning is necessary, but the choice of the bins can affect the recovered slopes. To test possible variations we repeat our analysis changing the bin size. Results are shown in Fig.~\\ref{fig:agefit}. Differences on the recovered slopes are within the errors and therefore in this case the age function is statistically not sensitive to the choice of binning.\n\nCaution must be taken when considering the age function at young ages: in the literature it has been proposed that an ``infant mortality'' (introduced by \\citealp{lada2003}), caused by the expulsion of leftover gas from star formation, could in principle cause a rapid decline in the number of clusters surviving after $\\sim10$ Myr. However numerical simulations show that gas expulsion does not have strong impact of the dynamical status of the stars within a gravitationally bound cluster (see \\citealp{longmore2014} review). As already discussed in \\citet{legus2}, at young ages it could be easier to include in the sample sources that are unbound at the origin. We are not considering the sizes of clusters, or their internal dynamics, therefore we are unable to assess the boundness of clusters. However, given a typical size of a few parsecs for the cluster radius \\citep{ryon2015,ryon2017}, we know that clusters older than 10 Myr have ages larger than their crossing time, and we can consider them bound systems. \nThe inclusion of unbound sources would cause a strong decline in the age function because they contaminate our sample at ages younger than 10 Myr. If we exclude the youngest age bins below 10 Myr, we find that the fit results in a shallower slope, $\\gamma=0.30 \\pm0.06$ (compare the resulting slopes in Fig.~\\ref{fig:agefit} and Fig~\\ref{fig:agefit_mass}). In the hypothesis of constant SFR over the last 200 Myr, we conclude that disruption is not very significant in the age range $10-200$ Myr. \n\nIn order to test if the disruption time is mass dependent, we divided the catalogue in three subsamples of increasing mass, in the ranges: $5000-10^4$ M$_\\odot$, $10^4-3\\times10^4$ M$_\\odot$\\ and masses $>3\\times10^4$ M$_\\odot$. The slopes of the resulting age functions (Fig.~\\ref{fig:agefit_mass}) present small differences, with more massive sources having flatter slopes. Incompleteness at ages $\\approx200$ Myr may start affecting the less massive sources ($\\sim$5000 M$_\\odot$) which seem to have significantly more disruption. The differences between the two more massive bins are within 1 $\\sigma$, thus very similar.\nCH16 found slopes compatible within 1 $\\sigma$ in different mass bins, i.e. $\\gamma=-0.71 \\pm0.03$, $-0.64 \\pm0.20$ and $-0.62 \\pm0.07$ for mass ranges log(M\/M$_\\odot$)$=3.8-4.0$ , $4.0-4.5$ and $4.5-5.0$ respectively. These slopes are systematically steeper than what we find, with differences close to $\\sim3 \\sigma$ from our values.\n\\begin{figure}\n \\centering\n \\includegraphics[width=\\columnwidth]{14_plots_agefunc_mass.pdf}\n\\caption{Age function divided in mass bins. Dashed lines represent the best fitting curve for the bins in the range Log(age\/yr)$=7-8.5$. The grey-shaded areas mark the part of the functions excluded from the analysis due to incompleteness (old ages) and possible contamination by unbound sources (young ages).}\n\\label{fig:agefit_mass}\n\\end{figure}\n\nThe difficulty in retrieving the correct model for mass disruption can be also due to the simultaneous action of different processes dispersing the clusters mass.\n\\citet{elmegreen2010} propose a model in which clusters are put into an hierarchical environment in both space and time and show that, under reasonable assumptions, many different processes of mass disruption (or the combination of them) can reproduce an age function with a power law decline, as generally observed.\n\nUnder the assumption of a MDD time, we can derive a typical value for cluster disruption. We consider $t_4$, i.e. the time necessary to disrupt a cluster of $10^4$ M$_\\odot$, as it is an indicative physical value. We use a maximum-likelihood fitting technique, introduced by \\citet{gieles2009}, where we assume an ICMF described by a $-2$ power law with a possible exponential truncation at $M_*$ (which is left as a free parameter) and a disruption process which is mass dependent in time, with a timescale-mass relation given by $t_{dis}\\propto M^{0.65}$. In this analysis we considered sources with ages up to $10^9$ yr, limiting the sample to $M>5000$ M$_\\odot$\\ and $V_{mag}<23.4$ mag (see Fig.~\\ref{fig:agefit_maxlk}). The mass cut, as previously pointed out, allows us the study of a mass-complete sample (up to 200 Myr). The cut in magnitude instead allows us to consider sources older than $200$ Myr, accounting for the fading of old sources below the magnitude completeness limit.\n\nThe results of this maximum-likelihood analysis are $M_*=(8.6\\ \\pm\\ 0.5)\\times10^4$ M$_\\odot$\\ and $t_4=230\\ \\pm\\ 20$ Myr (probability distributions are given in Fig.~\\ref{fig:agefit_maxlk}). \nThe Schechter truncation mass is compatible to what we found in the previous section within the uncertainties. The disruption timescale is of the same order of $t_4\\simeq130$ Myr found in the analysis of \\citet{gieles2009}, slightly above the interval $100\\lesssim t_4\\lesssim200$ Myr retrieved for different assumptions for the cluster formation history in \\citet{gieles2005}. However, the analysis of the age function suggests that some disruption, possibly also with mass independent time, may have been effective from the beginning, as we see a mild decrease in $dN\/dt$ already at young ages.\n\\begin{figure}\n \\centering\n\\subfigure{ \\includegraphics[width=0.95\\columnwidth]{15a_maxlk_age6to9_mcut37_magl23_4_tot_2ndpart.png}}\n\\subfigure{ \\includegraphics[width=1.\\columnwidth]{15b_maxlk_age6to9_mcut37_magl23_4_tot_1stpart.png}}\n\\caption{Probability distributions for the Schechter mass $M_*$ and for the typical disruption time of a $10^4$ M$_\\odot$\\ cluster, $t_4$, in the maximum-likelihood test (top). The red dot indicates the best fit. The plots on the bottom collect the age-mass distributions, with a dashed line indicating the limit (both due to the imposed mass cut and the limiting magnitude) above which we have selected the sources for the analysis. Black points are the selected sources, while blue points are the total sample. \n}\n\\label{fig:agefit_maxlk}\n\\end{figure}\n\n\n\n\\section{Simulated Montecarlo Populations}\n\\label{sec:montecarlo}\nIn order to better understand the results of the previous sections, we perform simulations, producing Monte Carlo populations of synthetic clusters, with properties similar to what we observe for the cluster population of M51. We use different initial sets of parameters for the mass distribution and the disruption and compare the synthetic populations with the observed one. This is particularly useful for the luminosity function, which, while easy to obtain, is the result of many generations of YSCs formed and evolved within the galaxy.\nAnalytical and semi-analytical models have tried to derive an expected LF shape from basic assumptions on the CIMF and on the age distribution and have been applied to the studies of cluster populations (e.g. \\citealp{haas2008} in M51, \\citealp{fall2006} in the Antennae, \\citealp{hollyhead2016} in NGC 1566). We will also refer to those studies in order to compare expectations, simulations and observations.\n\n\\subsection{Simulated Luminosity Functions}\nWe simulate populations of clusters with masses larger than 200 M$_\\odot$. The number of clusters per simulation is set such that we have the same number of clusters with $M>5000$ M$_\\odot$\\ as the observed population, i.e. $\\sim1200$. The star (and cluster) formation history is assumed constant for 200 Myr, which is also the maximum age we assign to the simulated clusters. After simulating age and mass for each of the synthetic clusters we assign a magnitude to them in the $B$, $V$ and $I$ bands, using the same models adopted for the SED fitting and described in Section~\\ref{sec:sed}. We can then build the luminosity functions, both in the binned and in the cumulative way (see Sec~\\ref{sec:lumfunc} for ref.), and fit them with a power law. The models used for the initial mass function and for the disruption are:\n\\begin{description}\n\\item[\\textbf{PL-2:}] pure $-2$ power law, no disruption.\n\\item[\\textbf{PL-2\\_MDD:}] pure $-2$ power law, mass dependent disruption time model with $k=0.65$ and $t_4=230$ Myr (as the results of the maximum-likelihood fit of Sec~\\ref{sec:agefunc} suggest, see the same Section also for the formalism).\n\\item[\\textbf{SCH:}] Schechter function\\footnote{The Schechter mass function, when not specified, is assumed to have a slope $-2$ in the power law part, i.e. $dN\\propto M^{-2}\\ e^{-M\/M_*} dM$} with $M_*=10^5$ M$_\\odot$, no disruption.\n\\item[\\textbf{SCH\\_MDD:}] Schechter function with $M_*=10^5$ M$_\\odot$, mass dependent disruption time model with $k=0.65$ and $t_4=230$ Myr.\n\\end{description}\nWe have not included MID in the models because it will only change the normalisation of the LF, not the shape. For this reason the two models without disruption (PL-2 and SCH) can be used to also study the expected LF in the case of mass independent disruption time model.\n\nThe results of the analysis are collected in Tab~\\ref{tab:mc_lum} and plotted in Fig~\\ref{fig:lumfit_simulated}. \nAs expected, a $-2$ power law mass function has a luminosity function with the same shape, when no disruption is considered. \nSlopes close to $-2$ are retrieved in all filters with both a binned and a cumulative function fit. The values for the cumulative function are slightly steeper, but both methods gives comparable results.\n\\begin{table}\n\\centering\n\\caption{Results of the fit of the luminosity functions derived from the simulated cluster populations with a pure power law. The models are described in the text. The fit of the observed luminosity function (with only sources younger than 200 Myr) is given for comparison. The errors, not reported in the table, are on the order of $\\pm0.03$ mag. The magnitude limits used are the same of the real data, listed in Tab.~\\ref{tab:lumfit}.}\n\\begin{tabular}{llllllll}\n\\hline\n\\multicolumn{1}{l}{MF}\t& \\multicolumn{3}{l}{binned function} \t& \\multicolumn{1}{c}{\\ }\t& \\multicolumn{3}{c}{cumulative function} \\\\\n\\\t\t\t\t\t& B \t\t& V \t\t& I\t\t\t\t& \\ \t\t\t\t\t& B \t\t& V \t\t& I \t\t\t\t\\\\\n\\hline\n\\hline\nOBS.\t\t\t\t& 1.94\t& 1.97\t& 1.99\t\t\t& \\\t\t\t\t\t& 2.14\t& 2.14\t& 2.25\t\t\t\\\\\n\\hline\nPL-2\t\t\t\t\t& 1.90\t& 2.00\t& 1.96\t\t\t& \\\t\t\t\t\t& 2.02\t& 2.03\t& 2.06\t\t\t\\\\\nPL-2\\_MDD\t\t\t& 1.65\t& 1.70\t& 1.76\t\t\t& \\\t\t\t\t\t& 1.78\t& 1.79\t& 1.83\t\t\t\\\\\nSCH\t\t\t\t\t& 1.94\t& 2.06\t& 2.13\t\t\t& \\\t\t\t\t\t& 2.19\t& 2.20\t& 2.38\t\t\t\\\\\nSCH\\_MDD\t\t\t& 1.81\t& 1.90\t& 2.02\t\t\t& \\\t\t\t\t\t& 2.08\t& 2.09\t& 2.29\t\t\t\\\\\n\\hline\n\\end{tabular}\n\\label{tab:mc_lum}\n\\end{table}\n\\begin{figure*}\n\\centering\n\\includegraphics[width=0.9\\textwidth]{16_lumfunc_simulated.pdf}\n\\caption{Luminosity functions of the simulated populations in the binned (left) and cumulative (right) form. $BVI$ bands are plotted in both cases. The models plotted are $-2$ power law (PL-2, solid line), $-2$ power law including MDD (PL-2\\_MDD, dashed) and Schechter (SCH, dashed-dotted) functions. Even if the underlying mass functions are very different, the luminosity functions are graphically quite similar in the binned form. When plotted cumulatively, instead, the differences are clear. In this second case the SCH model has the most similar LF shape to the one observed in the data.}\n\\label{fig:lumfit_simulated}\n\\end{figure*}\n\nConsidering disruption, MID would not change the shape of the mass or luminosity function, as discussed in Section~\\ref{sec:massage}. On the other hand a MDD would remove more quickly low-mass sources, modifying the luminosity function. The recovered slopes in this case are shallower, indicating that many low-luminosity sources have been removed (or have fallen below the completeness limit). \nThe effect of MDD is producing shallower LF in all filters. \nWe also observe that the slope is steeper in redder filters, as was also observed in the real data. This trend is mainly due to the difference in the the magnitude range fitted. The $I$ band, having a brighter completeness limit, is fitted only up to a magnitude of 22.25 mag, which is less affected by disruption than less bright magnitudes. If all filters will be fitted up to the same limiting magnitude the trend would not appear. \n\nWhen considering a Schechter mass function, the corresponding LF also has a steeper end. This is reflected in the recovered slopes, which in many cases have values more negative than $-2$. Including MDD still produces the effect of bending the low-luminosity end of the LF, but in this case, with slopes shallower than $-2$ at low luminosity and steeper than $-2$ at high luminosities, the resulting slope with a single power law can still be more negative than $-2$, as observed in the real data.\n\nWe can conclude that in order to produce luminosity functions with slopes steeper than $-2$ we need the underlying mass function to be truncated, or at least steeper than $-2$. MDD can affect the luminosity function, but only producing shallower functions. \n\n\n\n\\subsection{Age-Luminosity relation}\nAnother characteristic of the luminosity function that can be tested is the relative contribution of sources of different ages to the total luminosity in each magnitude bin. It has been proposed in the literature that, for a truncated mass function, the median age of clusters vary as function of luminosity, with the brightest clusters being on average younger than the faintest \\citep[see e.g. ][]{adamobastian2015}. This expectation has been confirmed with semi-analytical models \\citep[e.g.][]{larsen2009,gieles2010} and compared successfully with observations \\citep{larsen2009,bastian2012b}. \n\nTo verify that a similar trend is visible also with our catalogue we use the simulated cluster population, considering the PL-2 and SCH runs. \nFor both cases we divide the sources in magnitude bins of 1 mag and we take the median age of the sources inside each bin. We repeated this process 100 times and in Fig.~\\ref{fig:lumfit_ages} we plot the area covered by the distribution of the central 50\\% median ages per magnitude bin. In the case of PL-2 the median ages are more or less constant at all magnitude, even if towards the bright end the spread between the percentiles increases, due to the lower number of sources there. On the other hand, for SCH, a trend with younger ages towards brighter bins is clear. \nWe also include the median and 25th and 75th percent intervals obtained from the observed luminosities of the clusters. They show the same decreasing trend as the SCH model, bringing additional support to it. Similar conclusions are reached studying the age-luminosity relation for $\\sim$50 of the brightest clusters of M51 with spectroscopically-derived ages in a forthcoming paper (Cabrera-Ziri et al., in prep).\n\\begin{figure*}\n\\centering\n\\includegraphics[width=1\\textwidth]{17_plot_lumages_cut.pdf}\n\\caption{Median ages of the clusters in magnitude bins of 1 mag. The observed data (black points, with bars extending from the 1st to the 3rd quartiles) are plotted over the expectations from the models. The shaded areas encompass the central 50\\% distributions of medians from the simulations. The light areas (with an almost uniform behavior) are from the PL-2 model, the (declining) darker areas are from the SCH model. The data show a trend for younger ages associated in brighter clusters. This is in agreement with expectations from a Schechter mass function with a mass truncation at $10^5$ M$_\\odot$.}\n\\label{fig:lumfit_ages}\n\\end{figure*}\n\nWe have not considered disruption in this simple comparison. Anyway we do not expect the disruption to change drastically the results: MID is unable to produce the observed trend, and MDD could in principle only produce an opposite trend, with brighter sources being on average older \\citep{larsen2009}. We must conclude that the trend we see between ages and luminosities is another sign of an underlying truncated mass function.\n\n\n\\subsection{Simulated Mass Function}\n\\label{sec:mcmass}\nWe compare the simulated mass functions with the observed one. For each model we set the number of simulated clusters in order to be the same of the observed ones. In doing so we are able to test for the effect of random sampling from the mass function, which could produce a ``truncation-like'' effect (see \\citealp{dasilva2012}). We repeat the simulations 1000 times and compare the observed mass function with the median, the 50\\% and the 90\\% limits of the simulated functions in Fig.~\\ref{fig:massfit_simulated}. \nWe plot the mass functions in the cumulative form as we have seen that in this way the differences are graphically easier to spot. The models of the mass function considered for the simulations are a simple $-2$ power law, the single power law best fit of the cumulative function and a Schechter function with truncation mass $1.0\\times10^5$ M$_\\odot$\\ (best fit found in Tab.~\\ref{tab:massfit2}).\nWhen comparing the mass functions behaviours in Fig.~\\ref{fig:massfit_simulated} we immediately notice that a simple $-2$ power law (left panel) overestimates the number of cluster at high masses. A steeper power law (middle panel) follows better the observed data on average, but it underestimates the low mass clusters and overestimates the high mass ones. \nThe Schechter case (right panel) instead follows quite nicely the observed mass function at all masses.\n\\begin{figure*}\n\\centering\n\\subfigure{\\includegraphics[width=0.325\\textwidth]{18a_montecarlo_mass_pl2.pdf}}\n\\subfigure{\\includegraphics[width=0.325\\textwidth]{18b_montecarlo_mass_pl230.pdf}}\n\\subfigure{\\includegraphics[width=0.325\\textwidth]{18c_montecarlo_mass_sch.pdf}}\n\\caption{Monte Carlo simulations (black lines) compared to the observed mass function (blue triangles). 1000 Monte Carlo populations were simulated. The median mass distributions (solid lines) and the limits within 50\\% (dashed) and 90\\% (dotted) of the simulations are plotted. The resulting probability of the AD test is also reported (see text for description). The models for the simulated clusters are pure power laws with different slopes (left and center) and Schechter mass functions (right).}\n\\label{fig:massfit_simulated}\n\\end{figure*}\n\nWe test the null hypothesis that the observed masses are described by our models. We are mainly interested in the upper part of the mass function, which is the only part that possibly deviates from a simple power law description. We run the Anderson-Darling (AD) test comparing the observed masses larger than $1.0\\times10^4$ M$_\\odot$\\ with the ones produced in the Monte Carlo simulations. The AD test returns the probability that the null hypothesis (the two tested samples are drawn from the same distribution) is true and a typical value for rejecting the null hypothesis is $p\\sim10^{-4}$. The resulting probabilities of our test are collected in the plots of Fig~\\ref{fig:massfit_simulated}. \nThey confirm that $-2$ power law is a poor description for the massive part of the function (p $\\approx10^{-5}$) while the other cases perform better (p $=0.070$ for the steeper power law and p $=0.334$ for the Schechter function). Both for the steeper power law and for the Schechter function the null hypothesis is not rejected. The test gives a better agreement with the Schechter function.\n\nWe can conclude that the analysis of the mass function suggests that there is a mechanism that inhibits the formation of clusters at very high masses. As pointed out in Section~\\ref{sec:GMC}, the same is seen for GMCs, and this mass cut could therefore come from the progenitor structures.\n\n\n\n\\section{Cluster Formation Efficiency}\n\\label{sec:cfe}\nThe Cluster Formation Efficiency, CFE (also called $\\Gamma$), is the fraction of star formation happening within bound clusters \\citep[see][]{bastian2012}. Previous studies have found that $\\Gamma$ varies positively with the SFR density, $\\Sigma_{\\textrm{SFR}}$. Galaxies with higher $\\Sigma_{\\textrm{SFR}}$\\ also have on average a larger $\\Gamma$ (see \\citealp{goddard2010,adamo2011,ryon2014} and Fig. \\ref{fig:gamma}). Variations of CFE values have been observed also inside single galaxies, e.g. in the study of M83 \\citep{silvavilla2013,adamo2015}, suggesting that the main dependence of the efficiency is on the $\\Sigma_{\\textrm{SFR}}$.\n\nIn order to calculate $\\Gamma$, we compared the SFR with the cluster formation rate (CFR). We calculated the SFR from FUV, correcting for dust using 24$\\mu$m emission according to the recipe from \\citet{hao2011} and assuming a $0.1-100$ M$_\\odot$\\ Kroupa IMF. The resulting value is collected in Tab.~\\ref{tab:gamma}.\n\nThe CFR has been calculated in the age ranges, $1-10$, $1-100$, $1-200$ and $10-100$ Myr.\nThe sum of the masses of clusters in our catalogue with $M>5000$ M$_\\odot$\\ in the selected age range, divided by the age spanned, gives the CFR of the mass limited sample of sources. In order to correct for the missing mass of clusters less massive than 5000 M$_\\odot$, we have assumed a model for the ICMF and derived the contribution of the low-mass clusters down to 100 M$_\\odot$.\nThe assumed ICMF is a Schechter function with a truncation mass of $1.0\\times10^5$ M$_\\odot$, as resulting from the fit in Section~\\ref{sec:massfunc}.\n\nValues for measured masses, CFRs and $\\Gamma$ are collected in Tab.~\\ref{tab:gamma}. \nErrors on $\\Gamma$ have been estimated considering uncertainties of 0.1 dex on age and mass values, a Poisson error \non the number of clusters used for calculating the total mass, the errors on the mass function given in Tab.~\\ref{tab:massfit2} and assuming an uncertainty of $10\\%$ on the SFR value. \n$\\Gamma_{1-10}$ has been derived using two different tracers for the SFR. In one case H$\\alpha$+24$\\mu$m has been used following the recipe by \\citet{kennicutt2009}. This recipe assumes that SFR has been constant for $\\sim100$ Myr. In the second case H$\\alpha$\\ has been used without correction for obscured SF. This second case traces only the SF younger than 10 Myr but a 50\\% uncertainty is associated to the SFR value.\n\\begin{table}\n\\centering\n\\caption{CFR and CFE ($\\Gamma$) values for M51. $^a$ SFR derived from FUV+24$\\mu$m, associated uncertainty of $\\sim10\\%$. $^b$ SFR derived from H$\\alpha$+24$\\mu$m, associated uncertainty of $\\sim10\\%$. $^c$ SFR derived from H$\\alpha$\\ only, not corrected for obscured SF, associated uncertainty of $\\sim50\\%$.}\n\\begin{tabular}{ccccc}\n\\hline\nAge range\t\t& CFR\t\t& SFR\t\t& $\\langle\\Sigma_{\\textrm{SFR}}\\rangle$ \t\t\t\t& $\\Gamma$\t\\\\ \n\\ [Myr]\t\t& [M$_\\odot$\/yr]\t& [M$_\\odot$\/yr]\t& $\\left[\\textrm{M}_{\\odot}\/\\textrm{yr kpc}^{-2}\\right]$\t\t& [\\%]\t\\\\\n\\hline\n\\hline\n$10-100$\t\t& 0.305\t\t& 1.636$^a$\t\t\t& 0.0139\t\t& 18.6$\\ _{\\pm2.4}$\t\t\\\\\n\\hline\n$1-10$\t\t& 0.465\t\t& 1.437$^b$\t\t\t& 0.0132\t\t& 32.4$\\ _{\\pm12.1}$\t\\\\\n$1-10$\t\t& 0.465\t\t& 0.734$^c$\t\t\t& 0.0062\t\t& 63.3$\\ _{\\pm39.0}$\t\\\\\n$1-100$\t\t& 0.321\t\t& 1.636$^a$\t\t\t& 0.0139\t\t& 19.6$\\ _{\\pm2.5}$\t\t\\\\\n$1-200$\t\t& 0.264\t\t& 1.636$^a$\t\t\t& 0.0139\t\t& 16.2$\\ _{\\pm1.9}$\t\t\\\\\n\\hline\n\\end{tabular}\n\\label{tab:gamma}\n\\end{table}\n\n$\\Gamma_{10-100}=(18.6 \\pm2.4)\\%$ is probably the value that is less affected by systematics. The absence of clusters younger than 10 Myr effectively avoid the possible inclusion of unbound sources in the sample, while the restriction to ages younger than 100 Myr lowers the effects of the cluster disruption on the $\\Gamma$ derivation \\citep[see][]{kruijssen2016}.\nHowever, Tab.~\\ref{tab:gamma} shows that all $\\Gamma$ are consistent within 2$\\sigma$ with a 20\\% value. $\\Gamma_{1-10}$ and $\\Gamma_{1-100}$ are bigger due to the possible contamination by unbound sources already discussed in Section~\\ref{sec:agefunc}. $\\Gamma_{1-200}$ is slightly lower, but we know from the age function analysis (Section~\\ref{sec:agefunc}) that disruption is affecting the cluster population.\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=\\columnwidth]{19_gamma_m51_yesmodel.pdf}\n\\caption{Cluster Formation Efficiency $\\Gamma$ in function of the average $\\Sigma_{\\textrm{SFR}}$ derived in the age range 10-100 Myr (blue square). Values for other galaxies \\citep[taken from][]{goddard2010,adamo2011,annibali2011,ryon2014,adamo2015,lim2015,johnson2016} are shown for comparison. The black solid line is the $\\Sigma_{\\textrm{SFR}}-\\Gamma$ model presented in \\citet{kruijssen2012} with a 3$\\sigma$ uncertainty enclosed by the dotted lines.}\n\\label{fig:gamma}\n\\end{center}\n\\end{figure}\n\n$\\Gamma_{10-100}$ is compared with CFEs from other galaxies in Fig.~\\ref{fig:gamma}. Our value for M51 is similar to the $\\Gamma$ values found for other local galaxies, in particular for similar spiral galaxies like M83 \\citep{adamo2015}. More in general, it fits well into the $\\Sigma_{\\textrm{SFR}}$$-\\Gamma$ relation modeled by \\citet{kruijssen2012}, which predicts the amount of star formation happening in the cluster to increase with increasing surface density of star formation. \nIn a recent work on M 31, \\citet{johnson2016} suggested a comparison of the $\\Gamma$ values with the surface density of molecular gas. Their results show that the cluster formation efficiency scales with $\\Sigma_{\\textrm{gas}}$, with the mid-plane pressure of the galactic disk and with the fraction of molecular over atomic gas. These finding suggests that environments with higher pressure can form denser gas clouds which in turn result in a higher fraction of star formation happening in a clustered form. The environmental analysis of Paper II, based on regions with different $\\Sigma_{\\textrm{SFR}}$\\ and $\\Sigma_{\\textrm{gas}}$ inside M51, will help testing the cluster formation efficiency scenario.\n\n\n\n\n\n\n\n\n\\section{Conclusions}\n\\label{sec:conclusions}\nUsing LEGUS \\citep{legus1} broadband observations of M51 we have built a new catalogue of young stellar clusters. The new WFC3 coverage allows very accurate photometry in the $UV$ and $U$ bands, necessary for deriving precise ages of the young sources. \nThe cluster catalogue is automatically extracted using the steps described in \\citet{legus2}. A critical parameter of the extraction is the minimum concentration index considered (1.35 in our analysis). The SEDs of the extracted sources are fitted with \\textit{Yggdrasil} SSP models \\citep{yggdrasil} in a deterministic approach. \n\nSources with detection in at least 4 filters are initially visually classified by humans. This subsample has been used as a training set for a ML algorithm that has classified the entire catalogue.\nWe focus our analyses on 2834 sources which are compact and uniform in color, neglecting multi-peaked sources.\n\nClusters are spatially associated with the spiral arms of the galaxy up to ages of 200 Myr, as recently shown in a simulation of spiral galaxies by \\citet{dobbs2017}.\nThe luminosity, mass and age functions of the cluster sample are analyzed and compared to the results of simulated Monte Carlo cluster populations, in order to be able to better interpret the observed features. \n\nWe list hereafter the main results of this work.\n\\begin{itemize}\n\\item A double power law provides the best fit for the luminosity function, suggesting a truncation at bright magnitudes which is better observed when the function is plotted in a cumulative way. A trend of steeper slopes with redder filter is observed, as already pointed out by \\citet{haas2008}, and, in each filter, brighter sources have on average younger ages. \n\\item The mass function has been directly studied with a maximum--likelihood fit, supporting the hypothesis of a truncated function: the cluster population of M51 is consistent with a Schechter function with slope $-2$ and a mass truncation at $M_{*} = 10^5$ M$_\\odot$. \nThe analysis is repeated considering only high mass clusters (M$>10^4$ M$_\\odot$). A steeper slope is retrieved but the fit still gives preference to a truncated function with $M_{*}\\sim10^5$. \n\\item The age function indicates the presence of a moderate disruption over the range $10-200$ Myr.\nThe study of disruption in the first 10 Myr is precluded by the contamination of possibly unbound sources. Under the assumption of a mass-dependent disruption in time, a typical timescale for the disruption of $10^4$ M$_\\odot$\\ clusters is derived, $t_4=230$ Myr. \n\\item Simulated Monte Carlo populations are used to test the luminosity and mass function analyses. The trends observed in the luminosity function are recovered when a cluster population with an underlying Schechter mass function is analysed. A cluster population simulated with a pure power law mass function fails instead to produce the observed trends. \nWe notice that a cumulative function is the preferred way to study the bright end of the luminosity function and to investigate a possible deviation from a simple power law. The inclusion of cluster disruption in the simulations have also an impact on the luminosity function, but only at the faint end.\nMonte Carlo population are also used for a direct study of the mass function. A careful study of the high mass end of the function in the comparison with the simulated populations rejects the hypothesis of the function following a simple $-2$ power law, with a probability of $\\simeq0.001\\%$ given by an Anderson-Darling statistics. \n\\item We derive the fraction of star formation happening in clusters, which for M51 is $\\simeq20\\%$. This value is in line with the model of \\citet{kruijssen2012} linking the cluster formation efficiency with the SFR density, as well as consistent with observations for $\\gamma$ within local star-forming galaxies. \n\\end{itemize}\n\nIn a forthcoming paper (Paper II) clusters will be analyzed as a function of environments inside M51, enabling us to link cluster properties with the interstellar medium and GMC properties. \n\n\n\n\n\n\n\n\n\\section*{Acknowledgements}\nBased on observations made with the NASA\/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, \nwhich is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5--26555. \nThese observations are associated with program \\# 13364.\nSupport for program \\# 13364 was provided by NASA through a grant from the Space Telescope Science Institute.\nWe acknowledge the help and suggestions provided by M. Gieles and N. Bastian.\nCLD acknowledges funding from the European Research Council for the FP7 ERC starting grant project LOCALSTAR. \nDAG kindly acknowledges financial support by the German Research Foundation (DFG) through program GO 1659\/3-2.\nMF acknowledges support by the Science and Technology Facilities Council \n[grant number ST\/P000541\/1].\n\n\n\n\n\n\n\\bibliographystyle{mnras.bst}\n","meta":{"redpajama_set_name":"RedPajamaArXiv"}} +{"text":"\\section{Introduction}\n\nThe puzzling nature of the hidden order (HO) phase of URu$_2$Si$_2$\\ is still not understood. The central interest of this enigma is that it reflects the duality between the local and itinerant characters of the $5f$ electrons. These different facets often play a major role in the field of strongly correlated electronic systems. In spite of more than two decades of intense search\\cite{Palstra:1985}, there is still no direct access to the order parameter (OP) as it occurs for the sublattice magnetization of the high pressure antiferromagnetic (AF) ground state with the wave-vector \\textbf{Q$_{AF}$}=(0,0,1)\\cite{Amitsuka:1999,Bourdarot:2005} which appears above $P_x\\simeq0.5$ GPa via a first order transition switching from HO to AF phases\\cite{Motoyama:2003,Amitsuka:2007,Hassinger:2008}. However, recently it was noticed that an unambiguous signature of the HO phase is the sharp resonance at $E_0\\simeq1.8$ meV for the commensurate wave-vector \\textbf{Q$_0$}=(1,0,0) (equivalent wave-vector to \\textbf{Q$_{AF}$}) as this resonance mode collapses through $P_x$ while the other resonance at $E_1\\simeq4.1$ meV for the incommensurate wave-vector \\textbf{Q$_1$}=(1.4,0,0) persists through $P_x$\\cite{villaume:2008}. Furthermore, above $P_x$, a magnetic field leads to the ''resurrection'' of the resonance at $E_0$, when the HO phase is restored \\cite{aoki:2009}. As the strong inelastic signal at \\textbf{Q$_0$} is replaced above $P_x$ by a large elastic signal, fingerprint of the AF ground state with \\textbf{Q$_{AF}$}=(0,0,1), it was proposed that, in both HO and AF phases, a lattice doubling along the \\textbf{c} axis occurs at the transition from paramagnetic (PM) to either HO or AF ground states\\cite{aoki:2009}. The nice feedback is that the change in the class of tetragonal symmetry via development of a new Brillouin zone generates a drastic decrease in the carrier number as pointed out by a large number of theories\\cite{Yamagami:2000,Harima:2009,Elgazzar:2009} and experiments\\cite{Hassinger:2008a}. Furthermore, as the HO-AF line touches the PM-HO and PM-AF lines at a critical pressure $P_c\\simeq 1.4$ GPa and critical temperature $T_c\\simeq19.5$ K, a supplementary symmetry breaking must occur between the HO and AF phases. Two recent theoretical proposals for the OP of the HO phase were a hexadecapole - from Dynamical Mean Field Theory (DMFT) calculations\\cite{Haule:2009}; or a $O_{xy}$ type antiferroquadrupole - from group theory analysis \\cite{Harima::2010}. For both models, the pressure-switch from HO to AF will add a supplementary time reversal breaking in the AF phase. For these models as well as for recent band structure calculations \\cite{Elgazzar:2009}, partial gapping of the Fermi surface may happen at $T_0=17.8$ K with a characteristic gap $\\Delta_G$. It was even proposed in the last model, that in the HO phase the gapping is produced by a spontaneous symmetry breaking occurring through collective AF moment excitations.\n\nThis article presents a careful revisit of the inelastic neutron response at \\textbf{Q$_0$} \\cite{Broholm:1991,Mason:1995,Wiebe:2007} using recent progress in polarized inelastic neutron configuration of the spectrometer IN22 and in the performance of the cold-neutron three-axis spectrometer IN12, both installed at the high flux reactor of the Institute Laue Langevin (ILL). Thus, this new generation of experiments provide a careful basis on the temperature and pressure evolution of $E_0$. They give a new insight on previous neutron studies thanks to our recent proof\\cite{villaume:2008} that the resonance at $E_0$ is up to now the major signature of the HO phase. Comparison will be made with the case of charge ordering observed in the skutterudite PrRu$_4$P$_{12}$ as well as with PrFe$_4$P$_{12}$ where a sequence of HO-AF phases are observed \\cite{JPSJ77SA:2008}.\n\n\\section{Experimental set up}\n\nHigh quality single crystals of URu$_2$Si$_2$\\ were grown by the Czochralski method in tetra-arc furnace. The details are described elsewhere \\cite{aoki:2010}. The sample already used for inelastic neutron scattering in the superconducting state \\cite{Hassinger:2010}, was installed in an ILL-type orange cryostat with $\\bf{c}$ axis oriented vertically on the thermal triple-axis IN22 in its polarized configuration and on the cold triple-axis IN12 in its standard configuration (both CEA-CRG spectrometer at ILL).\n\nThe IN22 experiment was performed at 1.5 K, with two fixed final energies : 14.7 meV ($k_f$=2.662 \\AA$^{-1}$) and 30.3 meV ($k_f$=3.84 \\AA$^{-1}$). The beam was polarized by a Heusler monochromator vertically focusing and analyzed in energy and polarization by a Heusler analyzer vertically (fixed) and horizontally focusing. The flipping ratio was around 17 and the energy resolutions were 0.95 meV and 2.4 meV for both energies, respectively. No collimation was installed. The background was optimized by an optical calculation of the dimension openings of the slits placed before and after the sample. The measurements were performed in the non-spin-flip channel which provides all the necessary information. The inelastic scans were performed with $\\bf{Q}$ parallel to the $\\bf{a}$ axis. When the polarization is along the $\\bf{a}$ axis, the intensity measured (I$^a_{NSF}$) corresponds to the nuclear and background contributions; when the polarization is along the $\\bf{b}$ or $\\bf{c}$ axis, the intensities measured (I$^b_{NSF}$) and (I$^c_{NSF}$) correspond to the same contributions plus the (imaginary part of the) susceptibility along the $\\bf{b}$ or $\\bf{c}$ directions. In the following, the data shown are the subtraction of inelastic scans performed in the non-spin-flip channel with polarization along the $\\bf{b}$ or $\\bf{c}$ axis, with an inelastic scan performed in the same conditions with polarization along the $\\bf{a}$ axis. The subtractions correspond to the imaginary part of the dynamical susceptibilities with for the $\\bf{b}$ axis $\\chi''_y\\propto (I^b_{NSF}-I^a_{NSF})$ the transverse susceptibility and for the $\\bf{c}$ axis $\\chi''_z\\propto (I^c_{NSF}-I^a_{NSF})$ the longitudinal susceptibility.\n\nThe IN12 experiment was performed with a fixed final energy $E_f=4.7$ meV ($k_f$=1.5 \\AA$^{-1}$) that gives a good compromise between intensity of the excitation and energy resolution $l=0.11$ meV. The incident neutrons were selected by a (0,0,2) graphite vertically focusing monochromator with vertically focusing and analyzed in energy by a (0,0,2) graphite analyzer with horizontal focusing. No collimation was installed. The temperature was measured by a calibrated carbon thermometer and it was checked before each scan that the temperature was stable. As for IN22, the background was optimized by an optical calculation of the dimension openings of the slits placed before and after the sample. The raw scans were corrected for the electronic background (29 counts per hour) and for the $\\lambda\/2$ contamination of the monitor.\n\n\\section{Description of the model used for fitting the inelastic spectrum}\nIn a neutron scattering experiment, the neutron intensity $I(q,\\omega)$ in the detector is proportional to the convolution of the scattering function $S(q,\\omega)$ with the instrumental resolution function. $S(q,\\omega)$ is related to the imaginary part of the dynamical spin susceptibility $\\chi''(q,\\omega)$ ($\\chi(q,\\omega)= \\chi'(q,\\omega)+\\imath\\chi''(q,\\omega)$) via the fluctuation-dissipation theorem : $S(q,\\omega)=n(\\omega,T)\\chi''(q,\\omega)$ where $n(\\omega,T)= 1\/(1-e^{-\\hbar\\omega\/{k_BT}})$ is the detailed balance factor. \n Below $T_{0}$, the excitation is well-defined with an asymmetrical shape related to the finite extent of the resolution function that integrates the dispersion of the excitation in the vicinity of the nominal $q$ wave-vector at which the spectrometer is set-up. To analyze the data, firstly a harmonic oscillator function is taken for $\\chi\"(q,\\omega$) ; this corresponds to the difference of two normalized Lorentzian functions multiplied by $\\frac{\\chi_o\\,\\Omega^2_0}{\\omega_0\\gamma}$, where $\\chi_0=\\chi(q,\\omega=0)$ is the static susceptibility, $\\Omega_0$ is the oscillator frequency, $\\gamma$ its damping, and $\\omega_0$ is given by the equation $\\Omega_0=\\sqrt{\\omega^2_0+(\\gamma\/2)^2}$. Secondly, a simplified convolution with the resolution function is made : to this aim the resolution function is approximated by a 4D parallelepiped (instead of an ellipsoid) and the $q$ dispersion is taken as linear in all directions. This linear dispersion simplified the calculation and gives a better description of the dispersion than a usual quadratic law in the case of URu$_2$Si$_2$. This description is called the \\bm{$\\gamma$} \\textbf{model} when $\\gamma \\gg l$ and the \\bm{$l$} \\textbf{model} when $l \\gg \\gamma$, $l$ being the energy resolution. It leads to a simple analytic expression for $I(q,\\omega)$ at the minimums of the dispersion ($q_0$):\n\n\\begin{equation}\nI(q_0,\\omega)=n(\\omega,T)\\left(F(q_0,\\omega,)-F(q_0,-\\omega)\\right)\\\\\n\\label{m3}\n\\end{equation}\nwith\n\\begin{eqnarray}\n\\begin{split}\nF(q_0,\\omega)&=L\\,K(q_0,\\omega)\\left(\\sum_{i=x,y,z}\\sqrt{1+(2\\alpha_i\\,l\/\\omega_0(q_0))^2}\\right.\\nonumber\\\\\n&\\left.e^{-\\left(4\\ln{2}\\left(\\alpha_i\\frac{\\omega-\\omega_0(q_0)}{\\omega_0(q_0)}\\right)^2\\right)} \\right)\\\\\nK(q_0,\\omega)&=\\frac{1}{2}\\left(1+\\frac{2}{\\pi}arctan\\left(\\frac{2\\,(\\omega-\\omega_0(q_0))}{\\beta}\\right)\\right)\\nonumber\\\\\n\\end{split}\n\\end{eqnarray}\n\\\\\n\n\\noindent\nwhere $\\beta=\\gamma(q_0)$ for the \\bm{$\\gamma$} \\textbf{model} and $\\beta=l$ for the \\bm{$l$} \\textbf{model} \n$\\alpha_i=\\alpha_{i0}(\\omega_0\/\\Omega_0)^2$, where $\\alpha_{i0}$ is the ratio between the slope of the dispersion and the $q$-width of the resolution function. The L parameter, which depends of the magnetic form factor, the incident and final energies ($k_i$ and $k_f$) is supposed to stay constant at first approximation.\n\nFinally the susceptibility $\\chi_{0}$ for $TT_0$ (Fig. \\ref{fig12}), the magnetic response is over-damped, and the spectrum is treated with a quasi-elastic magnetic function of width $\\Gamma_L$ added to the continuum ($\\Gamma_c=7.8$ meV). As expected, $\\Gamma_L=\\gamma_0\/2$ at $T_0$, which validates the \\bm{$\\gamma$} \\textbf{model} and indicates that the gap $E_0$ drops to zero in the paramagnetic state (above $T_0$). The half-width of this excitation is plotted versus temperature in Figure \\ref{fig14} (open circles). $\\Gamma_L$ increases rapidly (maybe linearly) with temperature. It was not possible to follow this signal for temperatures approaching $T=27.5$ K: above this temperature only the large tail of the continuum can be detected.\n\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=80mm]{F12.eps}\n\\end{center}\n\\caption{Energy scans measured at \\textbf{Q$_0$} for temperatures just around the transition temperature $T_0=17.8$ K (and at 1.5 K for reference). Fits are described in the text. At $T=27. 1$ K (triangles), the scan is only fit by the continuum (quasi-elastic signal with $\\Gamma_c=7.8$ meV).}\n\\label{fig12}\n\\end{figure}\n\nAs the temperature evolution of the signal is not usual, the Figure \\ref{fig20} gives the variation of $\\chi''(\\mathbf{Q}_0,\\omega)\/\\omega$ as a function of $\\omega$ at different temperatures. Thus, this plot shows that this signal saturates at low temperature and the abrupt drop of E$_0$ when approaching $T_0$. Clearly the integration of $\\chi''(\\mathbf{Q}_0,\\omega)\/\\omega$ increases on cooling below $T_0$ as discussed latter as a consequence of Fermi Surface reconstruction.\n\nFigure \\ref{fig15} shows the magnetic susceptibility at $Q_0$, $\\chi(\\mathbf{Q}_0,\\omega=0)$. The susceptibility for $TT_0$ (open circles). Of course, in addition, the susceptibility coming from the magnetic continuum (with $\\Gamma_c=7.8$ meV) has to be added (constant low contribution in Fig \\ref{fig15}). Without any scaling-factor, the susceptibilities from the \\bm{$\\gamma$} \\textbf{model} and the second magnetic quasi-elastic contribution ($\\Gamma_L$) are equal at $T_0$. The total susceptibility versus temperature shows a saturation at low temperature, then decreases but with a bump around $T_0$. The susceptibility stays almost constant above $T = 30$ K. In contrast to Mason's analysis no marked divergence of $\\chi(\\mathbf{Q}_0,\\omega=0)$ is observed at $T_0$. Furthermore, in our data below 16 K, we found that $\\chi(\\mathbf{Q}_0,\\omega=0)$ increases on cooling before saturating at low temperature.\n\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=80mm]{F10.eps}\n\\end{center}\n\\caption{Temperature dependence of the gap $E_0$. The curves are guide for the eyes.}\n\\label{fig10}\n\\end{figure}\n\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=80mm]{F14.eps}\n\\end{center}\n\\caption{Temperature dependence of the half-width $\\gamma_0$ of the resonance at \\textbf{$Q_0$} below $T_0$ (filled circle) and of the quasi-elastic $\\Gamma_L$ above $T_0$ (open circle). The curve below $T_0$ corresponds to a fit with a Korringa model.}\n\\label{fig14}\n\\end{figure}\n\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=80mm]{F20.eps}\n\\end{center}\n\\caption{Imaginary part of the dynamical spin susceptibility without the magnetic continuum divided by energy $(\\chi''(\\mathbf{Q}_0,\\omega)\/\\omega$) for different temperatures. The curves are deduced from the $\\gamma$ model corrected by the thermal factor and divided by $\\omega$. The integration of these curves gives $\\chi(\\textbf{Q$_0$},\\omega=0)$.}\n\\label{fig20}\n\\end{figure}\n\nThe integration $\\int^{6.3 meV}_0\\chi''(\\mathbf{Q}_0,\\omega)d\\omega$ ($\\simeq I_{E_0}$) of the magnetic resonance at \\textbf{$Q_0$} without the continuum contribution ($\\Gamma_c$) is plotted in Fig. \\ref{fig13}. The filled circles correspond to the contribution with the oscillator model and the open circles correspond to the integration of the quasi-elastic contribution of width $\\Gamma_L$. This integration, below $T_0$, seems to mimic the temperature variation of an OP vanishing at $T_0$, while above $T_0$ it behaves like the contribution of a critical regime. Furthermore, the temperature variation of $I_{E_0}(T)$ is well described below $T_0$ by BCS formula used for the temperature variation of the superconducting gap \\cite{Bardeen:1957}.\n\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=80mm]{F15.eps}\n\\end{center}\n\\caption{Temperature dependence of the static susceptibility $\\chi(\\textbf{Q}_0,\\omega=0)$. In filled circles are data coming from the \\bm{$\\gamma$} \\textbf{model}, the open circles from the quasi-elastic model $\\Gamma_L$, and the hatched area coming from the magnetic continuum ($\\Gamma_c=7.8$ meV).}\n\\label{fig15}\n\\end{figure}\n\n\\begin{figure}\n\\begin{center}\n\\includegraphics[width=80mm]{F13.eps}\n\\end{center}\n\\caption{Temperature dependence of the integrated imaginary part of the dynamical spin susceptibility. In filled circles are data coming from the \\bm{$\\gamma$} \\textbf{model}, the open circles from the quasi-elastic model. The intensity coming from the magnetic continuum ($\\Gamma_c=7.8$ meV) was not taken into account. The line is the BCS-type gap below $T_0$, and a guide for the eyes above $T_0$.}\n\\label{fig13}\n\\end{figure}\n\n\\section {Discussion}\n\nThe occurrence of a magnetic continuum ($\\Gamma_c = 7.8$ meV) for energies $\\omega$ far higher than the energy $k_BT_0$ of the hidden order phase is a general phenomenon: large energy transfer leads to excite energy states not only from the ground state.\n\nAs the study is focused on the hidden order phase, we have concentrated our studies on the frequency response at \\textbf{Q$_0$}. Below $T_0$, the resonance at $E_0$ represents a collective mode with a half-width $\\gamma_0\/2$ much smaller than the gap energy $E_0$: the collective mode is long lifetime as shown by the large ratio $\\frac{E_0}{\\gamma_0\/2}\\approx35$. As shown on Fig. \\ref{fig10}, in decreasing temperature $E_0$ very rapidly reaches its final value $\\approx1.7$ meV. On crossing $T_0$, the resonance collapses and the magnetic response appears quasi-elastic with a half-width $\\Gamma_L$ becoming rapidly larger than the resonance energy $E_0$ at $T = 0$ K.\n\nThe susceptibility $\\chi($\\textbf{Q$_0$},$\\,\\omega$=0) reported in Fig. \\ref{fig15} is calculated for $T