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# Enhancements of Multi-class Support Vector Machine Construction from Binary
Learners using Generalization Performance
Patoomsiri Songsiri [email protected] Thimaporn Phetkaew
[email protected] Boonserm Kijsirikul [email protected] Department of
Computer Engineering, Chulalongkorn University, Pathumwan, Bangkok Thailand.
10330 School of Informatics, Walailak University, Thasala District, Nakhon Si
Thammarat Thailand. 80161
###### Abstract
We propose several novel methods for enhancing the multi-class SVMs by
applying the generalization performance of binary classifiers as the core
idea. This concept will be applied on the existing algorithms, i.e., the
Decision Directed Acyclic Graph (DDAG), the Adaptive Directed Acyclic Graphs
(ADAG), and Max Wins. Although in the previous approaches there have been many
attempts to use some information such as the margin size and the number of
support vectors as performance estimators for binary SVMs, they may not
accurately reflect the actual performance of the binary SVMs. We show that the
generalization ability evaluated via a cross-validation mechanism is more
suitable to directly extract the actual performance of binary SVMs. Our
methods are built around this performance measure, and each of them is crafted
to overcome the weakness of the previous algorithm. The proposed methods
include the Reordering Adaptive Directed Acyclic Graph (RADAG), Strong
Elimination of the classifiers (SE), Weak Elimination of the classifiers (WE),
and Voting based Candidate Filtering (VCF). Experimental results demonstrate
that our methods give significantly higher accuracy than all of the
traditional ones. Especially, WE provides significantly superior results
compared to Max Wins which is recognized as the state of the art algorithm in
terms of both accuracy and classification speed with two times faster in
average.
###### keywords:
support vector machine , multi-class classification , generalization
performance
††journal: Neurocomputing
## 1 Introduction
The support vector machine (SVM) [1, 2] is a high performance learning
algorithm constructing a hyperplane to separate two-class data by maximizing
the margin between them. There are two approaches for extending SVMs to multi-
class problems, i.e., solving the problem by formulating all classes of data
under a single optimization, and combining several two-class subproblems.
However, the difficulty and complexity to solve the problem with the first
method are due to the increase of the number of classes and the size of
training data, so the second method is more suitable for practical use. In
this paper, we focus on the second approach.
For constructing a multi-class classifier from binary ones, the method called
one-against-one trains each binary classifier on only two out of $N$ classes,
and builds $N(N-1)/2$ possible classifiers. Several strategies have been
proposed for combining the trained classifiers to make the final
classification for an unseen data. Friedman [3] suggested the combination
strategy called Max Wins. In the classification process of Max Wins, every
binary classifier provides one vote for its preferred class and the class with
the largest vote will be set to be the final output. Chang and Lee [4]
investigated an adaptive framework to manage a nuisance vote which is a vote
for an unrelated class by allowing a classifier to make a non-vote for data of
unrelated class. Instead of a binary classifier, they employed a ternary
classifier that consists of two particular classes and the rest of the classes
fused as the third class for solving this problem.
Vapnik [1] proposed the one-against-the-rest approach working by constructing
a set of $N$ binary classifiers in which each $i^{th}$ classifier is learned
from all examples in the $i^{th}$ class, and the remaining classes labeled
with the positive and negative classes, respectively. The class corresponding
to the classifier with the highest output value is used to make the final
output. Moreover, Manikandan and Venkataramani [5] adapted the traditional
one-against-the-rest to work as a sequential classifier. All classifiers will
be ordered corresponding to their misclassification. This method needs a lower
number of classifiers on avearge compared with the traditional one-against-
the-rest, but both algorithms have the same problem in the training phase
because of the difficulty for calculating the absolutely separating hyperplane
between a class and all of the other classes.
Dietterich and Bakiri [6] introduced the Error Correcting Output Code (ECOC)
based on the fundamental of information theory. For a given code matrix with
$N$ rows and $L$ columns, each element contains either $\textquoteleft$1’, or
$\textquoteleft$-1’. Each column denotes the bit string showing the
combination of positive and negative classes for constructing a binary
classifier, and each row of the code matrix indicates the unique bit string
for representing a specific class (each bit string is called a codeword).
Allwein et al. [7] extended the coding method by adding the third symbol
$\textquoteleft$0’ as “ don’t care bit” to allow the binary model learned
without considering some particular classes. Unlike the previous method, the
number of classes for training a binary classifier can be varied from 2 to $N$
classes. Based on these two systems for an $N$ classes problem, the maximum
numbers of different binary classifiers are $2^{N-1}-1$ [6], and
$\frac{3^{N}-2^{N+1}+1}{2}$ [8], respectively. Design of code matrices with
different subsets of binary classifiers gives different abilities for
separating classes, and the problem of selecting a suitable subset of binary
classifiers is complicated with a large size of $N$. To obtain the suitable
code matrix, some techniques using the Genetic Algorithm have been proposed
[9, 10]. In the classification phase, a test example is classified by all
classifiers corresponding to the column of the code matrix, and then the class
with the closet codeword is assigned to the final output class.
Platt et al. [11] proposed the Decision Directed Acyclic Graph (DDAG) in order
to reduce evaluation time [12]. In each round, a binary model will be randomly
selected from all $N(N-1)/2$ classifiers. The binary classification result is
employed to eliminate the candidate output classes, and to ignore all binary
classifiers related to the defeated class. It guarantees that the number of
classifications (applied classifiers) of the DDAG is always $N-1$. This
recursive task will be applied until there is only one last candidate class.
However, the misclassification of the DDAG can be occurred at the time when
selected binary classifiers related to the target class (hence forth BCRT)
give the wrong answer. The more times the number of BCRTs are applied, the
more chance the misclassification is produced by the DDAG. In order to reduce
this risk, Kijsirikul and Ussivakul [13] proposed the Adaptive Directed
Acyclic Graphs (ADAG) that has a reversed triangular structure of the DDAG. It
requires only $\lceil log_{2}N\rceil$ times or less that the target class is
tested against the other classes, while the DDAG possibly requires at most
$N-1$ times.
In addition, there have been many attempts that apply some information such as
the margin size [11], the number of support vectors [14], and the separability
measures between classes [15, 16], to improve the performance of the multi-
class classification. The margin size and the number fo support vectors were
applied for selecting the suitable two-class classifiers in the DDAG [11, 14].
The separability measure was employed for automatically constructing a binary
tree of multi-class classification based on the concept of the minimum
spanning tree [15]. Li, et al. [16] used similar information to vote the
preferred class for data in unclassifiable region for both the one-against-one
and the one-against-the-rest techniques.
In this research, we investigate the framework for enhancing three well-known
methods, which are the DDAG, the ADAG, and Max Wins. Max Wins is currently
recognized as the-state-of-the-art combining algorithm and it is also the most
powerful technique among all of our focused works with a need of $N(N-1)/2$
number of classifications for an $N$-class problem, while the other two
approaches reduce the number of classifications to $N-1$. We study the
characteristics of these methods that lead to wrong classification results.
The first two techniques have the same hierarchical structure and have the
same weak point that they “trust on individual opinion” for making decision to
discard the candidate classes. Intuitively, if only one of BCRTs makes a
mistake, the whole system will give the wrong output. The last technique as
the high performance one, Max Wins is based on the concept of “trust on most
popular opinion” for making decision to select the output class. If all of
$N-1$ BCRTs give the correct answer, Max Wins will always provide the correct
output class. However, if there exists only one of BCRTs give the wrong
answer, it may lead to misclassification due to equal voting, or other non-
target classes reaching the largest vote as shown later in the paper. Examples
which are incorrectly classified in this scenario can be recovered by our
proposed strategies.
In this paper, we demonstrate that the above traditional methods can be
improved based on the same idea that if we access further important
information of generalization performance of all binary classifiers and
properly estimate it, it can be employed for enhancing the performance of the
methods. Based on this idea, we propose four novel approaches including (1)
the Reordering Adaptive Directed Acyclic Graph (RADAG), (2) Strong Elimination
of the classifiers (SE), (3) Weak Elimination of the classifiers (WE), and (4)
Voting based Candidate Filtering (VCF). The first approach, the next two
approaches, and the last approach are improved from the ADAG, the DDAG, and
Max Wins, respectively. We also empirically evaluate our methods by comparing
them with the traditional methods on the sixteen datasets from the UCI Machine
Learning Repository [17].
This paper is organized as follows. Section 2 reviews the traditional multi-
class classification frameworks. Section 3 describes how to properly estimate
the generalization performance of binary classifiers. Section 4 presents our
proposed methodologies. Section 5 performs experiments and explains the
results and discussions. Section 6 concludes the research.
## 2 Multi-class Support Vector Machines
### 2.1 Max Wins
For an $N$-class problem, all possible pairs of two-class data are learned for
constructing $N(N-1)/2$ classifiers. All binary classifiers are applied for
voting the preferred class. A class with maximum vote will be assigned as the
final output class. This method is called Max Wins [3]. However, in case that
there exists more than one class giving the same maximum vote, the final
output class can be obtained by random selection from candidate classes with
the equal maximum-vote. An example of the classification using this technique
for a four-class problem is shown in Fig. 1. Each class will be voted (solid-
line) or ignored (dash-line) by all related binary models. For example, class
$1$ has three related classifiers, i.e., 1 vs 2, 1 vs 3, and 1 vs 4. The
voting result of class $1$, class $2$, class $3$ and class $4$ are three,
zero, one, and, two, respectively. In this case, class $1$ has the largest
score, and therefore it is assigned as the final output class.
Figure 1: An example of a four-class classification with Max Wins.
### 2.2 Decision Directed Acyclic Graphs
Platt et al. [11] introduced a learning algorithm using the Directed Acyclic
Graph (DAG) to represent the classification task, called the Decision Directed
Acyclic Graph (DDAG). This architecture represents a set of nodes connected by
edges with no cycles. Each edge has an orientation and each node has either 0
or 2 edges. Among these nodes, there exists a root node which is the unique
node with no edge pointing into it. In a DDAG, the nodes are arranged with a
triangular shape in which each node is labeled with an element of a boolean
function. There exists a single root node at the top, two nodes in the second
layer, and so on until the final layer of $N$ leaves for an $N$-class problem.
To make a classification, an example with an unknown class label is evaluated
by the nodes as binary decision functions. The binary output result in each
layer is applied to eliminate the candidate output classes and the binary
classifiers related to the defeated class are removed. At the first layer (see
Fig. 2), the root node can be randomly selected from all possible $N(N-1)/2$
classifiers and there are $N$ candidate output classes. After the root node is
tested, its binary result is employed to eliminate the candidate output
classes and the binary classifiers corresponding to the defeated class are
discarded. In the next layer, the remaining binary classifiers are randomly
selected to continue the same process in which some classes are eliminated
from the remaining candidate classes. The process is repeated until there is
only one class remained which is then assigned as the final output class. This
algorithm requires only $N-1$ decision nodes in order to obtain the final
answer.
Figure 2: The DDAG finding the best class out of four classes [11].
One disadvantage of the DDAG is that its classification result is affected by
the sequence of binary classifiers randomly selected in the evaluation
process. Platt et al. also proposed the other method that prefers the binary
decision function with the higher generalization performance measured by its
margin sizes, called the large margin DAGs [11]. The margin size ($\Delta$) is
a parameter for bounding the generalization ability of the binary SVM as shown
in terms of the VC-dimension in Eq. (2). It illustrates that the
generalization performance of the binary model is proportional to the size of
the margin. A binary classifier with the larger margin size will be firstly
applied in each round of the evaluation step. Moreover, Takahashi and Abe [14]
proposed a similar framework that employed the number of support vectors as a
performance measure. In this method, the generalization error
($\epsilon_{ij}$) for classes $i$ and $j$ was bounded by Eq. (1) [18]:
$\epsilon_{ij}=\frac{SV_{ij}}{M_{ij}},$ (1)
where $SV_{ij}$ is the number of support vectors for classes $i$ and $j$ and
$M_{ij}$ is the number of training data for classes $i$ and $j$.
### 2.3 Adaptive Directed Acyclic Graphs
Figure 3: The structure of an adaptive DAG for an 8-class problem.
In the DDAG, binary classification result of a previously employed binary
classifier is used to eliminate a candidate output classe, and there are only
current remaining candidate classes that can be possibly assigned as the final
output class. Therefore, the misclassification of a selected BCRT is the
crucial point.
The ADAG was originally designed to reduce this risk of the DDAG by using
reversed triangular structure [13]. In an $N$-class problem, there are
$\lceil\frac{N}{2}\rceil$ nodes at the top, $N/2^{2}$ nodes in the second
layer and so on until the lowest layer of the final node, as illustrated in
Fig. 3. Like the DDAG, binary output results of the ADAG in each layer are
applied to discard candidate output classes and the binary classifiers related
to the defeated classes are also ignored. Therefore, the ADAG also evaluates
only $N-1$ nodes to obtain the final answer.
According to the critical issue of misclassification mentioned above, even
only one selected classifier related to the target class provides a wrong
answer, the misclassification on the final output class cannot be avoided.
Hence, the number of times the target class is tested against other classes
indicates the risk of misclassification. The DDAG requires at most $N-1$ times
that the target class is tested against other classes, while the ADAG requires
only $\lceil log_{2}N\rceil$ times or less. This shows that the opportunities
of the target class tested against other classes on the ADAG is much lower
than the DDAG.
## 3 An Estimation of the Generalization Performan-ces of Binary Support
Vector Machines
The generalization performance of a learning model is the actual performance
evaluated on unseen data. For support vector machines, a model is trained by
using the concept of the Structure Risk Minimization principle [19] in which
the generalization performance of the model is estimated based on both terms
of the complexity of model (the VC dimension of approximating functions) and
the quality of fitting training data (empirical error). Consider the problem
of binary classification where dataset $X$ of $m$ samples in real
$n$-dimensional space is randomly independent identically distributed
observations drawn according to $P(x,y)=P(x)P(y|x)$. The expected risk
($R(\alpha)$) with probability at least $1-\delta$ can be bounded by the
following equation [20, 21]:
$R(\alpha)\leq\frac{l}{m}+\sqrt{\frac{c}{m}({\frac{R^{2}}{\Delta^{2}}{log^{2}m}+log{\frac{1}{\delta}}})},$
(2)
where there is a corresponding constant $c$ for all probability distributions,
$l$ is the number of labeled examples in $z$ with margin less than $\Delta$,
$R$ indicates the radius of the smallest sphere that contains all the data
points, and $\Delta$ is the distance between the hyperplane and the closest
points of the training set (margin size). The first and second terms of
inequality in Eq. (2) denote the bound of the empirical error, and the VC
dimension, respectively.
In our frameworks, the generalization ability will be applied to improve the
multi-class classification. Although there have been many attempts to use some
performance measures such as the margin size [11], the number of support
vectors [14], they may not accurately reflect the actual performance of each
binary SVM. Consider a two-class problem where hyperplanes $h_{1}$ and $h_{2}$
are learning models created to separate the positive and the negative
examples. Suppose that they provide different margin sizes of $\Delta_{1}$ and
$\Delta_{2}$, and the different numbers, $l_{1}$ and $l_{2}$, of labeled
examples in $z$ with the margin less than their margin sizes, respectively. In
case that the parameters $c$ and $\delta$ are fixed, there are only two
parameters including $\Delta$ and $l$ that affect the performance of the
learning model (as the parameters $m_{1}$ and $m_{2}$, as well as $R_{1}$ and
$R_{2}$ are the same for the same pair of a two-class problem). Now consider
two-learning models learned from different pairs of a two-class problem. In
case that the parameters $c$, and $\delta$ are fixed, according to inequality
in Eq. (2), obviously, if we use only $\Delta$, $\l$, or combination of them,
they are not sufficient to represent the whole term of their generalization
abilities. This shows that a binary model with the larger margin size may not
provide more accurate result of classification. The use of only the number of
support vectors is also shown in [21] that it is not predictive for
generalization ability.
(a)
(b)
(c)
(d)
Figure 4: Generalization errors of 325 classifiers of the Letter dataset based
on $k$-fold cross-validation (CV Bound), the number of support vectors (shown
in term of the ratio between the number of support vectors and the number of
training data: SV Bound), the margin size (shown in term of its inverse value
normalized to be in [0,1]: Normalized Margin Bound), and their actual risks on
test data (unseen data) by applying the polynomial kernel of $d=4$. Figure (a)
compares generalization errors calculated by all techniques where classifiers
are sorted in the ascending order by their actual generalization performances
(actual risk), and figures (b)-(d) show the comparisons between the actual
risks, and the estimated generalization errors with different measures, i.e.,
CV Bound, SV Bound, and Normalized Margin Bound, respectively (the classifiers
will be sorted in ascending order by the estimated generalization errors, and
for ease of visualization we show only $10\%$ of classifiers by sampling every
ten classifers from the sorted list of the classifiers.
Figure 5: Classification process of the RADAG.
As described above, the generalization ability can be employed to enhance the
performance of multi-class classification, by carefully design algorithms
which utilize this information as a selection measure for good classifiers. We
believe that the generalization performance of binary SVMs can be directly
estimated by $k$-fold cross-validation [22] (see Algorithm 1), and it can be
used to fairly compare the performances of binary SVMs on different two-class
problems. Below we give an example which demonstrates that $k$-fold cross-
validation is more suitable for estimating the generalizaiton performance of
the classifiers than the other measures used by the previous methods, i.e. the
number of support vectors, the margin size.
Algorithm 1 An estimation of the generalization error of a classifier by using
$k$-fold cross-validation.
1:procedure Cross Validation
2: Set of training data $T$ is partitioned into $k$ disjoint equal-sized
subsets
3: Initial the classification error of round $i$: $\epsilon_{i}\leftarrow 0$
4: for $i$=1 to $k$ do
5: validate set $\leftarrow$ $i^{th}$ subset
6: training set $\leftarrow$ all remaining subsets
7: Learn model based on training set
8: $\epsilon_{i}$ $\leftarrow$ Evaluate the learned model by validate set, and
find the number of examples with misclassification
9: end for
10: generalization error $\leftarrow$
$\sum_{i=1}^{k}\epsilon_{i}\times\frac{1}{|T|}$
11: return generalization error
12:end procedure
Fig. 4 shows the generalization performance measured by the previous methods
[11, 14], and $k$-fold cross-validation, which we propose to use as the
performance measure, for the Letter dataset with 26 classes, by applying the
polynomial kernel of $d=4$. Fig. 4 (a) illustrates that the trend of estimated
generalization error by $k$-fold cross-validation is very closed to the actual
risk, while the other two techniques give high variation. To further
investigation in more details, we select about $10\%$ of all classifiers to
show in Fig. 4 (b-d); these figures illustrate the comparisons between the
actual risk and the estimated generalization errors with different measures,
i.e., CV Bound, SV Bound, and Normalized Margin Bound, respectively. In each
figure, classifiers are sorted in ascending order by the estimated
generalization errors. It is expected that if a specific measure is a good
estimator for generalization error, its value should be in the same trend as
the actual risk (its value shoud increase with the increase of the actual
risk). A good trend is found in Fig. 4 (b), while the other two methods give
no clear trend and contain confusing patterns. In order to evaluate the
efficiency of these estimating methods, we apply the correlation analysis
between two variables [23], i.e., the actual risks and these three estimated
generalization errors. These evaluations are based on 325 classifiers as in
Fig. 4 (a), and the statistical r-values of them are 0.805, 0.372, and -0.230
as shown in Fig. 4 (b-d), respectively. The r-values also confirm that CV
Bound and actual risk have high correlation, while the other two methods give
very low correlation. They show that $k$-fold cross-validation is more
suitable to be the measure for the performance of binary classifiers.
According to the above reason, we apply this measure in our research.
## 4 The proposed methods
The combination of binary SVMs with high generalization performance directly
affects the accuracy of the multi-class classification. In this section, we
introduce four enhanced approaches based on the previous techniques i.e., the
ADAG, the DDAG, and Max Wins by applying the generalization abilities in order
to select suitable binary classifiers. An improvement of the ADAG is called
the Reordering Adaptive Directed Acyclic Graph (RADAG). There are two improved
versions for the DDAG i.e., Strong Elimination of the classifiers (SE) and
Weak Elimination of the classifiers (WE). The last technique is Voting-based
Candidate Filtering (VCF) enhanced from Max Wins. To increase the
classification accuracy, the generalization estimated by $k$-fold cross-
validation is utilized as the goodness measure of classifiers in our
frameworks.
### 4.1 Reordering Adaptive Directed Acyclic Graph
The ADAG is designed to reduce the number of times the binary classifiers
related to the target class are applied, from at most $N-1$ times required by
the DDAG, to $\lceil log_{2}N\rceil$ times or less. However, binary
classifiers in the first level of the ADAG are still randomly selected, and
its misclassification can be produced at the time even when only one BCRT
gives a wrong answer. In this section, we introduce a more effective method
which uses the minimum weight perfect matching to select the optimal pairs of
classes in each level with minimum generalization error. We called the method
the Reordering Adaptive Directed Acyclic Graph (RADAG).
The structure of the RADAG is similar to the ADAG, but they are different in
the initialization of the binary classifiers in the top level and the order of
classes in lower levels (see Fig. 4.1). The reordering algorithm with minimum
weight perfect matching is described in Algorithm 2. The algorithm selects the
optimal order of classes in each level. It is different from the ADAG in that
the initial order of classes in the ADAG is obtained randomly, and the
matching of classes in successive levels depends on the classification results
of nodes from the previous level. For the RADAG, the reordering process will
be applied to the remaining candidate classes in all levels for determining
the optimal sequence of them.
Figure 6: (a) A graph for an 8-class problem (b) An example of the output of
the reordering algorithm.
To select the optimal set of classifiers, the generalization measure in
Section 3 is used as a criterion. This scheme provides less chance to predict
the wrong class from all possible
$\frac{N!}{2^{\lfloor{N/2}\rfloor}{\lfloor{N/2}\rfloor}!}$ orders. Among
$N(N-1)/2$ classifiers, $N/2$ classifiers which have the smallest sum of
generalization errors will be used in the classification.
Algorithm 2 Reordering Adaptive Directed Acyclic Graph (RADAG).
1:procedure RADAG
2: Initial set of candidate output classes $C=\\{1,2,3,...,N\\}$, and set of
discarded classes $D=\emptyset$
3: Calculate generalization errors of all possible pairs of classes on $C$ as
described in Section 3
4: Create the binary SVMs from all possible pairs of classes on $C$
5: while $|C|>1$ do
6: Apply the minimum weight perfect matching [24] to find the optimal
$\lfloor\frac{|C|}{2}\rfloor$ pairs of classes from all possible pairs on $C$
to obtain the optimal binary models with minimum generalization error
7: $D$ $\leftarrow$ Classify the example by the optimal binary models, and
find the defeated classes
8: $C$ $\leftarrow$ $C-D$
9: end while
10: final output class $\leftarrow$ the last remaining candidate class
11: return final output class
12:end procedure
Let $G=(V,E)$ be a graph with node set $V$ and edge set $E$. Each node in $G$
denotes one class and each edge indicates one binary classifier of which
generalization error is estimated from Section 3 (see Fig. 6(a)). The output
of the reordering algorithm for graph $G$ is a subset of edges with the
minimum sum of generalization errors of all edges and each node in $G$ is met
by exactly one edge in the subset (see Fig. 6(b)).
Given a real weight $\epsilon_{e}$ being generalization error for each edge
$e$ of $G$, the problem of reordering algorithm can be solved by the minimum
weight perfect matching [24] that finds a perfect matching $M$ of minimum
weight $\sum(\epsilon_{e}:e\in M)$.
For $U\subseteq V$, let $E(U)=\\{(i,j):(i,j)\in E,i\in U,j\in U\\}$. $E(U)$ is
the set of edges with both endpoints in $U$. The set of edges incident to node
$i$ in the node-edge incidence matrix is denoted by $\delta(i)$. The convex
hull of perfect matchings on a graph $G=(V,E)$ with $|V|$ even is given by
a) $x\in\\{0,1\\}^{m}$
b) $\sum_{e\in\delta(v)}x_{e}=1$ for $v\in V$
c) $\sum_{e\in E(U)}x_{e}\leq\lfloor\frac{|U|}{2}\rfloor$ for all odd sets
$U\subseteq V$ with $|U|\geq 3$ or by (a),(b) and
d) $\sum_{e\in\delta(U)}x_{e}\geq 1$ for all odd sets $U\subseteq V$ with
$|U|\geq 3$
where $|E|=m$, and $x_{e}=1$ ($x_{e}=0$) means that $e$ is (is not) in the
matching.
Hence, the minimum weight of a perfect matching is at least as large as the
value of
$min\sum_{e\in E}\epsilon_{e}x_{e}$ (3)
where $x$ satisfies “(a), (b), and (c)” or “(a), (b) and, (d)”. Therefore, the
reordering problem can be solved by the integer program in Eq. (3).
### 4.2 Strong & Weak Elimination of Classifiers for Enhancing Decision
Directed Acyclic Graph
Algorithm 3 Strong Elimination of the classifiers (SE).
1:procedure SE
2: Initial set of candidate output classes $C=\\{1,2,3,...,N\\}$, and set of
discarded classes $D=\emptyset$
3: Calculate generalization errors of all possible pairs of classes on $C$ as
described in Section 3
4: Create the binary models from all possible pairs of classes on $C$
5: Sort the list of the binary models in ascending order by the generalization
errors
6: current classifier $\leftarrow$ the first element on the sorted list
7: while $|C|>1$ do
8: $D$ $\leftarrow$ Classify the example by the current classifier, and find
the defeated class
9: $C$ $\leftarrow$ $C-D$
10: current classifier $\leftarrow$ the next element on the sorted list where
it is not related to any classes discarded from $C$
11: end while
12: final output class $\leftarrow$ the last remaining candidate class
13: return final output class
14:end procedure
Algorithm 4 Weak Elimination of the classifiers (WE).
1:procedure WE
2: Initial set of candidate output classes $C=\\{1,2,3,...,N\\}$, and set of
discarded classes $D=\emptyset$
3: Calculate generalization errors of all possible pairs of classes on $C$ as
elaborated in section 3
4: Create the binary models from all possible pairs of classes on $C$
5: Sort the list of the binary models in ascending order by the generalization
errors
6: current classifier $\leftarrow$ the first element on the sorted list
7: while $|C|>1$ do
8: $D$ $\leftarrow$ Classify the example by current classifier, and find the
defeated class
9: $C$ $\leftarrow$ $C-D$
10: current classifier $\leftarrow$ the next element on the sorted list where
it does not include all two classes discarded from $C$
11: end while
12: final output class $\leftarrow$ the last remaining candidate class
13: return final output class
14:end procedure
According to the characteristic of the DDAG, binary classification results of
the previously employed binary classifiers are used to eliminate the candidate
output classes, and thus the final output class will be assigned with one of
the remaining candidate classes. By using the random technique for selecting a
binary classifier, the DDAG produces mis-classification at the time when a
BCRT with very low performance is selected and provides the wrong answer, as
the target class will be discarded from the remaining candidate classes, and
it is not possible to reach the correct output class. In this section, we
propose the framework to enhance the performance of the DDAG to select the
binary classifier with high performance based on the generalization abilities
of binary classifiers as described in Section 3.
We propose two methods that are Strong Elimination of the classifiers (SE) and
Weak Elimination of the classifiers (WE). Both algorithms are described in
Algorithm 3 and Algorithm 4. We also show a classification process of SE and
WE for an $N$-class problem in Fig. 7 and Fig. 8, respectively.
For both of the DDAG and SE, in each round, a defeated class will be removed
from candidate output classes, and all binary classifiers related to the
defeated class are ignored. Due to this reason, they guarantee $N-1$ number of
classifications for an $N$-class problem. However, these ignored classifiers
may have high generalization abilities and thus are helpful to eliminate the
other remaining candidate classes. Therefore, we then propose WE to make use
of binary classifiers with high generalization abilities.
According to the classifier elimination of WE, the number of classifications
is bounded with the best case of $N-1$, and the worst case of $N(N-1)/2$.
However, WE provides the opportunities to employ better classifiers as shown
in the Fig. 8. At round $r$, suppose that classifier $A_{i}$ vs $A_{j}$ has
lower generalization error than classifier $A_{i}$ vs $A_{k}$ , and both of
them are active classifiers. In this case, it is possible that classifier
$A_{i}$ vs $A_{j}$ can remove the class $A_{i}$ from the list of two remaining
candidate classes, and can avoid using classifier $A_{i}$ vs $A_{k}$ with
lower reliability that is unavoidable for SE as shown in Fig. 7.
Figure 7: Classification process of SE for an $N$-class problem.
Figure 8: Classification process of WE for an $N$-class problem.
### 4.3 Voting Based Candidate Filtering
Max Wins is one of high performance techniques that work based on the concept
of “trust on the most popular opinion” for making decision to select the
output class. If all of $N-1$ BCRTs give the correct answer, Max Wins will
always provide the correct output class. It does not depend on the answers of
the other binary classifiers. However, if only one of BCRTs gives a wrong
answer, it may lead to misclassification due to equal voting, or another non-
target class reaching the largest vote. Fig. 9 shows an example of such cases,
taken from our experiment on the Letter dataset (see Section for more
details); Fig. 9(a) and (b) show the cases of equal voting and another non-
target class having the largest vote, respectively.
We propose a novel multi-class classification approach that alleviates the
above problem of Max Wins, and uses the same concept “trust on the most
popular opinion” for filtering out the low competitive classes. On the other
hand, high competitive classes will be voted to be candidate output classes,
though there exist some BCRTs providing the wrong answer. If there is more
than one remaining class, the output class will be selected via the mechanism
of WE. Our proposed technique aims to combine the strong point of both Max
Wins and WE, and is called Voting based Candidate Filtering (VCF). The details
of our algorithm are shown in Algorithm 5.
Algorithm 5 Voting based Candidate Filtering (VCF).
1:procedure VCF
2: Initial set of candidate output classes $C=\\{1,2,3,...,N\\}$, and score of
class $i$: $s_{i\in N}\leftarrow 0$
3: Create the binary models from all possible pairs of classes on $C$
4: for $j$=1 to $N(N-1)/2$ do
5: $w$ $\leftarrow$ Classify the example by classifier $j^{th}$, and find the
winner class
6: $s_{w}$ $\leftarrow$ $s_{w}+1$
7: end for
8: $s_{top}$ $\leftarrow$ Find the top voting score of all $s_{i\in N}$
9: for $i$=1 to $N$ do
10: $dp_{i}\leftarrow\frac{(s_{top}-s_{i})\times 100}{s_{top}}$
11: if $dp_{i}\leq$ $threshold\\_value$ then
12: Add class $i$ into the set of candidate output classes $C$
13: end if
14: end for
15: if $|C|>1$ then
16: final output class $\leftarrow$ Call the WE procedure
17: else
18: final output class $\leftarrow$ the last remaining candidate class
19: end if
20: return final output class
21:end procedure
Let $s_{top}$, and $s_{i}$ indicate the maximum of scores for all $N$ classes,
and the score of class $i\in[N]$ for a test data, respectively. Also let
$dp_{i}$ denotes the percentage of the difference between $s_{top}$ and
$s_{i}$. An example of the calculation of $dp_{i}$ is shown in Fig. 9 (a),
where $i=\textquoteleft$E’, the score of class $\textquoteleft$E’ = 23 points,
and the score of class $\textquoteleft$C’ = 24 points (as the top score). Then
$dp_{i}$ value can be calculated by $\frac{(24-23)\times 100}{24}=4.17$. We
also define $threshold\\_value$ to be the threshold of $dp_{i}$ for
considering class $i$ as a candidate for the target class; class $i$ will be
accepted into the set of high competitive candidate classes if and only if its
$dp_{i}$ is less than or equal to $threshold\\_value$. We want to keep the
size of the filtered candidate classes as small as possible while still
containing the target class.
(a)
(b)
Figure 9: An example of high risk of misclassification of Max Wins together
with score distribution of all classes: two cases of misclassification of
class $\textquoteleft$C’ due to only one BCRT giving the wrong answer in the
Letter problem having 26 classes (25 possible BCRTs, and 25 points as the
largest possible score), a) three classes, including $\textquoteleft$C’,
$\textquoteleft$G’, and $\textquoteleft$L’, with equal score (only one BCRT
$\textquoteleft$C vs G’ giving the wrong class), and b) the non-target class
$\textquoteleft$E’ with the highest score (only one BCRT $\textquoteleft$C vs
E’ providing the wrong class).
Figure 10: A case study of the $461$ examples with high risk of
misclassification in the Letter problem. The maximum voting scores of these
examples are reached (1) by both of the target class and the non$-$target
class (equal vote: $1^{st}$ rank), or (2) by a non$-$target class (absolutely
wrong: $2^{nd}-8^{th}$ rank). The figure shows the target class score of these
examples by observing between the $dp_{t}$ and the rank of the target class.
A case study of high risk of misclassification in the Letter dataset including
4,010 examples where Max Wins provides 3,549 examples with the correct result,
and 461 examples with high risk of misclassification. By a high-risk example,
we mean (1) the example with an equal vote (the score of the target class is
equal to those of other non-target classes) and (2) the example with a vote
less than the maximum vote that is then mis-classified by Max Wins. These high
risk examples will be hopefully recovered with the correct class label by our
proposed algorithm. In our experiment, the high-risk examples includes 24
examples (around 5$\%$) with an equal vote, and 437 examples (around 95$\%$)
with a vote less than the maximum as shown in Fig. 10, where $dp_{t}$
represents the percentage of the difference between $s_{top}$ and the score of
the target class. For each example, we calculate the rank of the voting score
of the target class compared to the other non-target classes, and consider
only the first eight ranks. There are 24 examples (around $5\%$) in the first
rank, while in the second to the eighth ranks, the numbers of examples are
171, 77, 31, 30, 15, 14, and 10 (around $39\%,18\%,7\%,7\%,3\%,3\%$, and
$2\%$), respectively. The examples with the different ranks have different
ranges of $dp_{t}$ values, such as, in the second rank, the $dp_{t}$ values
are varied from $4.0$ to $12.0$, in the third rank, the $dp_{t}$ values are
varied from $4.2$ to $20.0$, in the fourth rank, the $dp_{t}$ values are
varied from $8.7$ to $20.0$, and so on.
According to this case study, there can be at most $5\%$ of examples that will
be correctly classified with the correct class label by random selection of
Max Wins, while the other $95\%$ of examples will be absolutely misclassified.
We want to recover an example that is not correctly classified by Max Wins, as
its actual target class is not in the first rank or its target class has equal
vote with some other output classes. If $threshold\\_value$ is set as $1$ in
the VCF algorithm, it will guarantee that all high-risk misclassified examples
with $dp_{t}$ values no greater than can be filtered into the set of the
candidate output classes; in this case only the examples in the first rank
($5\%$ of examples) will be selected. When we apply a bigger threshold, e.g.
$threshold\\_value$ $=10$, it covers all misclassified examples in the first
and the second ranks ($5\%+39\%$), almost of the third rank ($18\%$), and some
parts of the fourth rank ($7\%$). It shows that the increase of
$threshold\\_value$ covers more candidate classes, while the larger size of
$threshold\\_value$ creates a higher risk to employ an unnecessarily large
number of binary classifiers. On the other hand, if $threshold\\_value$ is too
low, the target class may be removed. However, a suitable $threshold\\_value$
can be obtained by general tuning techniques. For our experiment, we just
define $threshold\\_value$ to be $10$ for all of datasets without fine-tuning
which is good enough to demonstrate the effectiveness of the VCF algorithm.
## 5 Experiments
In this section, we design the experimental setting to evaluate the
performance of the proposed methods. We compare our methods with the
traditional algorithms, i.e., the DDAG, the ADAG, and Max Wins. We divide this
section into two parts as experimental protocols, and results & discussions.
### 5.1 Experimental Protocol
We run experiments on sixteen datasets from the UCI Machine Learning
Repository [17] including Page Block, Glass, Segment, Arrhyth, Mfeat-factor,
Mfeat-fourier, Mfeat-karhunen, Mfeat-zernike, Optdigit, Pendigit, Primary
tumor, Libras Movement, Abalone, Krkopt, Spectrometer, and Letter (see Table
1). For the datasets containing both training data and test data, we added up
both of them into one set, and used 5-fold cross validation for evaluating the
classification accuracy.
Table 1: Description of the datasets used in the experiments. Datasets | #Cases | #Classes | #Features
---|---|---|---
Page Block | 5,473 | 5 | 10
Glass | 214 | 6 | 9
Segment | 2,310 | 7 | 18
Arrhyth | 438 | 9 | 255
Mfeat-factor | 2,000 | 10 | 216
Mfeat-fourier | 2,000 | 10 | 76
Mfeat-karhunen | 2,000 | 10 | 64
Mfeat-zernike | 2,000 | 10 | 47
Optdigit | 5,620 | 10 | 62
Pendigit | 10,992 | 10 | 16
Primary tumor | 315 | 13 | 15
Libras Movement | 360 | 15 | 90
Abalone | 4,098 | 16 | 8
Krkopt | 28,056 | 18 | 6
Spectrometer | 475 | 21 | 101
Letter | 20,052 | 26 | 16
In these experiments, we scaled data to be in [-1,1] and employed two kernel
functions i.e., the Polynomial kernel $K(x_{i},x_{j})\equiv|({\bf
x_{i}}\cdot{\bf x_{j}}+1)|^{d}$, and the RBF kernel $K(x_{i},x_{j})\equiv
e^{-\gamma||{\bf x_{i}}-{\bf x_{j}}||^{2}}$. For the polynomial kernel we
applied the same set of degrees $d=\\{2,3,4,5\\}$ to all datasets, and for the
RBF kernel we applied the set of degrees $\gamma_{1}=\\{1,0.5,0.1,0.05\\}$ to
Page Block, Glass, Segment, Mfeat-zernike, Pendigit, Libras Movement, Abalone,
Krkopt, and Letter, and applied the set of degrees
$\gamma_{2}=\\{0.1,0.05,0.01,\\\ 0.005\\}$ to the other datasets. The default
parameter of regularization parameter $C$ was used for model construction;
this parameter is used to trade off between error of the SVM on training data
and margin maximization. In the training phase, we used software package
$SVM^{light}$ version 6.02 [25, 26] to create the $N(N-1)/2$ binary
classifiers. For the DDAG and the ADAG, we examined all possible orders of
classes for datasets having not more than 8 classes, whereas we randomly
selected 50,000 orders for datasets having more than 8 classes, and we then
calculated the average of accuracy of these orders.
### 5.2 Results & Discussions
We compare the original methods with their enhanced techniques in three tasks
including: (1) the ADAG with the RADAG, (2) the DDAG with two improved
approaches, i.e., SE and WE, and (3) Max Wins with VCF.
We also selected the best techniques from (1) and (2), i.e., the RADAG and WE,
respectively, and compared them with Max Wins as the state of the art
technique. These comparison results are shown in Table 2 to Table 5. Moreover,
paired comparison among all of three traditional methods (the DDAG, the ADAG,
and Max Wins), and all proposed techniques (SE, the RADAG, WE, and VCF) are
concluded in Table 6.
The best accuracy among these methods is represented in bold-face. In
addition, we used the one-tailed paired t-test technique to analyze the
significant difference between the accuracies of the traditional algorithms
and the proposed algorithms. To estimate the difference between accuracies, we
use a $k$-fold cross-validation method [22].
To indicate the level of the confidence interval using a one-tailed paired
t-test in the Table 2 to Table 5, the symbol ‘$+$’ and ‘$-$’ are used to
represent that the corresponding method has higher accuracy, and lower
accuracy compared to a baseline method, respectively. The number of symbols
shows the level of confidence interval for estimating the difference between
accuracies of two algorithms i.e., one symbol, two symbols, and three symbols
represent 90%, 95%, and 99% respectively.
Table 2: A comparison of the classification accuracy of the ADAG and the RADAG. | Polynomial | | RBF
---|---|---|---
Data sets | ADAG | RADAG | | ADAG | RADAG
Page Block | 93.597 | 93.541 $---$ | | 93.562 | 93.555 $-$
Glass | 63.879 | 64.019 | | 63.084 | 63.318
Segment | 93.207 | 93.236$+{\color[rgb]{1,1,1}*}$ | | 93.348 | 93.366
Arrhyth | 63.489 | 63.470 | | 58.049 | 57.991
Mfeat-factor | 97.238 | 97.225 $--{\color[rgb]{1,1,1}*}$ | | 96.921 | 96.938
Mfeat-fourier | 82.839 | 82.863 | | 82.456 | 82.513$+$
Mfeat-karhunen | 96.864 | 96.863 | | 96.890 | 96.900
Mfeat-zernike | 82.368 | 82.413$+++{\color[rgb]{1,1,1}*}$ | | 81.867 | 81.888
Optdigit | 98.995 | 98.999 | | 98.620 | 98.630
Pendigit | 99.400 | 99.402 | | 99.313 | 99.320$+$
Primary tumor | 47.266 | 47.619$++{\color[rgb]{1,1,1}*}$ | | 46.089 | 46.429$+++$
Libras Movement | 73.218 | 73.194 | | 72.289 | 72.569$+$
Abalone | 27.603 | 27.648 | | 27.353 | 27.337
Krkopt | 53.102 | 53.239$+++$ | | 53.088 | 53.173$+++$
Spectrometer | 54.445 | 54.842$++{\color[rgb]{1,1,1}*}$ | | 50.808 | 51.579$+{\color[rgb]{1,1,1}*}$
Letter | 88.668 | 88.787$+++$ | | 89.989 | 90.090$+++$
The experimental results in Table 2 uses the ADAG as the baseline algorithm.
It shows that the RADAG yields highest accuracy in several datasets. The
results also show that, at $95$% confidence interval, the RADAG performs
statistically better than the ADAG in five datasets using the Polynomial
kernel and better in three datasets using the RBF kernel. As shown in the
table, the RADAG performs better when the number of classes is comparatively
large, and does not perform well in the datasets with the small number of
classes, i.e., the Page Block, and the Mfeat-factor with 5 and 10 classes,
respectively. We believe that in case of datasets with the large number of
classes, the variety of generalization errors of classifiers in consideration
is rich and the RADAG is able to choose good classifiers freely, whereas the
RADAG may be forced to select ineffective classifiers in case of the small
number of classes, and it could lead to an incorrect output class.
Table 3: A comparison of the classification accuracy between the DDAG, and our methods, i.e, SE, and WE. | Polynomial | | RBF
---|---|---|---
Data sets | DDAG | SE | WE | | DDAG | SE | WE
Page Block | 93.597 | | 93.541 $---$ | 93.623$+$ | | 93.562 | | 93.555 $-$ | 93.582$++{\color[rgb]{1,1,1}*}$
Glass | 63.892 | | 64.019 | 64.019 | | 63.084 | | 63.201 | 63.201
Segment | 93.207 | | 93.236 $+$ | 93.247$++$ | | 93.350 | | 93.344 | 93.366
Arrhyth | 63.490 | | 63.527 | 63.527 | | 58.048 | | 57.991 | 58.162
Mfeat-factor | 97.238 | | 97.250 | 97.238 | | 96.923 | | 96.975$+$ | 96.975$++{\color[rgb]{1,1,1}*}$
Mfeat-fourier | 82.837 | | 82.863 | 82.863 | | 82.443 | | 82.475 | 82.538
Mfeat-karhunen | 96.863 | | 96.875 | 96.875 | | 96.861 | | 96.850 | 96.988$+{\color[rgb]{1,1,1}*}$
Mfeat-zernike | 82.362 | | 82.400$++{\color[rgb]{1,1,1}*}$ | 82.350 | | 81.869 | | 81.888 | 81.863
Optdigit | 98.994 | | 99.013$++{\color[rgb]{1,1,1}*}$ | 99.008 $+{\color[rgb]{1,1,1}*}$ | | 98.618 | | 98.643$+$ | 98.630
Pendigit | 99.399 | | 99.404 | 99.402 | | 99.312 | | 99.318 | 99.320$+$
Primary tumor | 47.227 | | 47.064 | 47.460$++$ | | 46.019 | | 46.032 | 46.111
Libras Movement | 73.142 | | 73.264 | 73.472$++{\color[rgb]{1,1,1}*}$ | | 72.283 | | 72.569$+$ | 72.431
Abalone | 27.611 | | 27.648 | 27.672$+$ | | 27.354 | | 27.330 | 27.398
Krkopt | 53.101 | | 53.263 $++$ | 53.472$+++$ | | 53.088 | | 53.212 $+++$ | 53.320$+++$
Spectrometer | 54.373 | | 54.632 | 54.421 | | 50.821 | | 51.316 | 51.842$++$
Letter | 88.609 | | 88.707 $+++$ | 88.835$+++$ | | 89.903 | | 89.977 $++$ | 90.294$+++$
Table 3 shows the experimental results of SE and WE compared with the DDAG as
the baseline algorithm. Both WE and SE have higher accuracy than the
traditional DDAG in almost all datasets. The results also show that at $95$%
confidence interval, SE performs statistically significantly better than the
DDAG in four datasets using the Polynomial kernel and significantly better
than the DDAG in two datasets using the RBF kernel. It is similar to the
previous comparison between the ADAG and the RADAG that in datasets with the
small number of classes, the classifier manipulation of SE may be forced to
select inaccurate classifiers and it possibly leads to the misclassification.
The results also show that WE performs statistically significantly better than
the DDAG in five datasets in both cases of the Polynomial kernel and the RBF
kernel. These results illustrate that WE can reduce the risk of selecting
inaccurate classifiers compared to SE.
Figure 11: An example of generalization errors of binary SVMs used by WE and SE in the Letter dataset. Table 4: A comparison of the classification accuracy of Max Wins and VCF. Data sets | Polynomial | | RBF
---|---|---|---
| Max Wins | VCF | | Max wins | VCF
Page Block | 93.600 | 93.623 | | 93.567 | 93.582$++{\color[rgb]{1,1,1}*}$
Glass | 63.863 | 64.019 | | 63.143 | 63.201
Segment | 93.209 | 93.247$++{\color[rgb]{1,1,1}*}$ | | 93.351 | 93.366
Arrhyth | 63.489 | 63.527 | | 58.048 | 58.162
Mfeat-factor | 97.242 | 97.238 | | 96.927 | 96.975$++{\color[rgb]{1,1,1}*}$
Mfeat-fourier | 82.852 | 82.825 | | 82.454 | 82.525
Mfeat-karhunen | 96.879 | 96.875 | | 96.952 | 96.963
Mfeat-zernike | 82.338 | 82.350 | | 81.825 | 81.863
Optdigit | 99.004 | 99.013 | | 98.631 | 98.630
Pendigit | 99.402 | 99.402 | | 99.315 | 99.320
Primary tumor | 47.394 | 47.460 | | 46.508 | 46.191
Libras Movement | 73.194 | 73.472$++{\color[rgb]{1,1,1}*}$ | | 72.373 | 72.431
Abalone | 27.614 | 27.672$+$ | | 27.375 | 27.398
Krkopt | 53.149 | 53.475$+++$ | | 53.146 | 53.328$+++$
Spectrometer | 54.263 | 54.421$+{\color[rgb]{1,1,1}*}$ | | 51.026 | 51.842$++$
Letter | 88.706 | 88.869$+++$ | | 90.112 | 90.316$+++$
Table 5: A comparison of the classification accuracy between Max Wins and the RADAG, and WE. | Polynomial | | RBF
---|---|---|---
Data sets | Max Wins | RADAG | WE | | Max Wins | RADAG | WE
Page Block | 93.600 | 93.541 $---$ | 93.623 | | 93.567 | 93.555 $--{\color[rgb]{1,1,1}*}$ | 93.582$++{\color[rgb]{1,1,1}*}$
Glass | 63.863 | 64.019 | 64.019 | | 63.143 | 63.318 | 63.201
Segment | 93.209 | 93.236 $+{\color[rgb]{1,1,1}*}$ | 93.247$++$ | | 93.351 | 93.366 | 93.366
Arrhyth | 63.489 | 63.470 | 63.527 | | 58.048 | 57.991 | 58.162
Mfeat-factor | 97.242 | 97.225 $---$ | 97.238 | | 96.927 | 96.938 | 96.975$++{\color[rgb]{1,1,1}*}$
Mfeat-fourier | 82.852 | 82.863 | 82.863 | | 82.454 | 82.513 | 82.538
Mfeat-karhunen | 96.879 | 96.863 | 96.875 | | 96.952 | 96.900 | 96.988
Mfeat-zernike | 82.338 | 82.413$++{\color[rgb]{1,1,1}*}$ | 82.350 | | 81.825 | 81.888$++{\color[rgb]{1,1,1}*}$ | 81.863
Optdigit | 99.004 | 98.999 | 99.008 | | 98.631 | 98.630 | 98.630
Pendigit | 99.402 | 99.402 | 99.402 | | 99.315 | 99.320 | 99.320
Primary tumor | 47.394 | 47.619$++$ | 47.460 | | 46.508 | 46.429 | 46.111
Libras Movement | 73.194 | 73.194 | 73.472$++{\color[rgb]{1,1,1}*}$ | | 72.373 | 72.569 | 72.431
Abalone | 27.614 | 27.648 | 27.672$+$ | | 27.375 | 27.337 | 27.398
Krkopt | 53.149 | 53.239 $++{\color[rgb]{1,1,1}*}$ | 53.472$+++$ | | 53.146 | 53.173 | 53.320$+++$
Spectrometer | 54.263 | 54.842$++$ | 54.421 $+$ | | 51.026 | 51.579 | 51.842$++$
Letter | 88.706 | 88.787 $++{\color[rgb]{1,1,1}*}$ | 88.835$+++$ | | 90.112 | 90.090 | 90.294$+++$
We further analyze the results comparing WE and SE on the Letter dataset which
consists of $26$ classes and 325 binary learners, as shown in Fig. 11. These
325 classifiers in the figure are sorted in ascending order by the
generalization error, and this sequence of classifiers is maintained in the
classification phase. SE requires $25$ classifiers and WE requires $93$
classifiers in this case, and the generalization error of the worst binary
classifier in WE is almost five times lower than in SE (the largest
generalization errors of all binary SVMs used in SE and WE are 0.015, 0.073,
respectively). As a result, the average performance of the binary classifiers
in WE is higher than SE.
As shown in Table 4 with Max Wins as the baseline method, VCF yields higher
accuracy than Max Wins in almost all of datasets. The results show that, at
$95$% confidence interval, in the Polynomial kernel VCF performs statistically
significantly better than Max Wins in four datasets, and in the RBF kernel VCF
performs statistically significantly better than Max Wins in five datasets.
The previous three tables show that our proposed methods improve the accuracy
of the ADAG, the DDAG, and Max Wins significantly.
Next, we select the best algorithm in each table from the first two tables,
i.e, the RADAG, and WE, and then compare them to Max Wins. According to
experimental result in Table 5, at $95$% confidence interval, the RADAG
performs statistically significantly better than Max Wins in five datasets
using the Polynomial kernel, and significantly higher than Max Wins in one
dataset using the RBF kernel. In case of the small number of classes, it is
possible that the RADAG will have the effect mentioned above. For WE, the
results show that it performs statistically significantly better than Max Wins
in four datasets in case of the Polynomial kernel and significantly better
than Max Wins in five datasets in case of the RBF kernel. There is no any
dataset in which Max Wins has significantly higher accuracy than WE.
Table 6: Paired comparisons among all techniques including of three traditional techniques (DDAG, ADAG, and Max Wins), and four proposed techniques (RADAG, SE, WE, and VCF). | | Traditional Methods | | Proposed Methods
---|---|---|---|---
Kernel Function | Algorithms | DDAG | ADAG | Max Wins | | SE | RADAG | WE | VCF
Polynomial | DDAG | | 1-15-0 | 2-14-0 | | 4-11-1 | 5-9-2 | 5-11-0 | 6-10-0
| ADAG | | | 2-13-1 | | 4-11-1 | 5-9-2 | 4-12-0 | 5-11-0
| Max Wins | | | | | 3-12-1 | 5-9-2 | 4-12-0 | 4-12-0
| SE | | | | | | 2-14-0 | 4-11-1 | 4-11-1
| RADAG | | | | | | | 2-12-2 | 4-10-2
| WE | | | | | | | | 1-15-0
RBF | DDAG | | 2-14-0 | 2-14-0 | | 2-14-0 | 3-13-0 | 5-11-0 | 5-11-0
| ADAG | | | 2-14-0 | | 1-15-0 | 3-13-0 | 5-11-0 | 5-11-0
| Max Wins | | | | | 2-12-2 | 1-14-1 | 5-11-0 | 5-11-0
| SE | | | | | | 2-13-1 | 4-12-0 | 3-13-0
| RADAG | | | | | | | 3-13-0 | 3-13-0
| WE | | | | | | | | 1-15-0
Table 6 summarizes paired comparisons of all algorithms including the
traditional techniques, and the proposed works based on both of the Polynomial
kernel and the RBF kernel. We show the win-draw-loss record (s) of the
algorithm in the column against the algorithm in the row. A win-draw-loss
record reports how many datasets the method in the column is better than the
method in the row (win), is equal (draw), or is worse (loss) at $95$%
confidence interval. As summarized in the table, our proposed methods are
better than all previous works i.e., the DDAG, the ADAG, and Max Wins. WE and
VCF give the highest accuracy among all of our methods. The result also shows
that VCF gives a little better results compared to WE. However, as mentioned
before in Section 4.3, the accuracies of VCF are the ones without fine-tuning,
and higher accuracies can be expected if fine-tuning is performed to find the
optimal $threshold\\_value$ for VCF.
### 5.3 Computational Time
The computational times of all methods are shown in Fig. 12 and Fig. 13. We
can classify algorithms according to the time requirement into three groups,
for an $N$-class problem: 1) $N-1$ times i.e., the DDAG, the ADAG, SE, and the
RADAG, 2) average about half of time of $N(N-1)/2$ i.e., WE, 3) $N(N-1)/2$
i.e., Max Wins, and VCF.
The results show that algorithms in the first and the second groups require
comparatively low running time in all datasets, especially when the number of
classes is relatively large, while the larger the number of classes, the more
running time the algorithms in the third group requires. WE in the second
group requires $N-1$ classifiers in the best case and $N(N-1)/2$ classifiers
in the worst case; however, in our experimental results WE takes approximately
half of time required by the algorithms in the third group. For the RADAG,
though the number of classes affects the running time for reordering process,
it takes a little time even when there are many classes. The algorithms in the
third group need $O(N^{2})$ comparisons for a problem with $N$ classes. VCF
needs more time to choose the final class from the set of candidate classes
which can be obtained by re-using the previous results of binary
classification.
The DDAG reduces the number of comparisons down to $O(N)$. SE spends a little
time more than the DDAG for sorting the classifiers in the training phase. By
reducing the depth of the path, the ADAG and SE require $O(N)$ comparisons of
binary classifiers. WE consumes more time than SE due to each round of
classification can reduce only one classifier while SE can eliminate all
classifiers built from the discarded class. The number of testing classifiers
for WE is equal to that for Max Wins in the worst case; fortunately, the
experimental results show that WE actually spends only half of Max Wins’ times
in the average case. The RADAG needs a little time more than the ADAG for
reordering the order of classes. Note that, the minimum weight perfect
matching algorithm, which is used in the reordering algorithm, runs in time
bounded by $O(N(M+NlogN))$ [24], where $N$ is the number of nodes (classes) in
the graph and $M=N(N-1)/2$ is the number of edges (binary classifiers). The
RADAG will reorder the order of classes in every level, except for the last
level. The order of classes in the top level is reordered only once and we use
the order to evaluate every test example. Hence for classifying each test
data, we need $log_{2}N-2$ times of reordering, where each time the number of
classes is reduced by half. Therefore, the running time of the RADAG is
bounded by $O(c_{1}N)+O(c_{2}N^{3}log_{2}N)$, where $c_{1}$ is much larger
than $c_{2}$.
Figure 12: A comparison of the computational time using the Polynomial kernel.
Figure 13: A comparison of the computational time using the RBF kernel.
## 6 Conclusion
Max Wins is a powerful combining technique with a need of $N(N-1)/2$ number of
classifications for an $N$-class problem, while the DDAG and the ADAG reduce
the number of classifications to $N-1$. We study the characteristics of these
previous methods that lead to wrong classification results. We believe that
the performances of them depend on the BCRTs. In case of Max Wins, if there
exists only one BCRT giving an incorrect answer, it may convey
misclassification due to equal voting or another non-target class reaching the
largest vote, while in cases of the DDAG and the ADAG, if only one of BCRTs in
the sequence of selected classifiers makes a mistake, the whole system will
give the wrong output. We investigate the well-organized combination of the
binary models including BCRTs in classification process to provide a more
precise final result.
In this research, we propose four methods for overcoming the above weakness of
the previous works. All our proposed methods are based on the same principle
that if the information about genearalization ability is accurately measured,
then it is able to be employed for enhancing the performance of the
classification. In this paper, the generalization performance is estimated by
$k$-fold cross-validation technique, and we show that it is more suitable than
previously used measures in other frameworks, such as the margin size and the
number of support vectors. Our proposed methods are the Reordering Adaptive
Directed Acyclic Graph (RADAG), Strong Elimination of the classifiers (SE),
Weak Elimination of the classifiers (WE), and Voting based Candidate (VCF).
The RADAG is an enhanced version for the ADAG by using the minimum weight
perfect matching for selecting the optimal pair of classes in each level with
minimum generalization error. Compared to the ADAG, the RADAG is not only
superior in terms of accuracy, but also maintains the same testing time
($N-1$). Next, We propose two improved algorithms for the DDAG, i.e. SE and
WE. In SE, a sequence of binary classifiers selected by minimum generalization
error is applied to eliminate the candidate classes until only one class
remained and assigned as the final output class. SE provides better accuracy
than the DDAG. The testing time of SE is the same as the traditional DDAG and
the RADAG, with a number of applied classifiers equal to $N-1$. We also
propose the other enhanced version for the DDAG, called WE. This approach aims
to efficiently use as many as possible of the classifiers with low
generalization errors. This is different from the process of the DDAG and SE
in which all binary classifiers related to a defeated class are ignored when
the defeated class is removed from the candidate classes. In WE, however, a
classifier will be ignored only if all of two related classed of that
classifier are discarded from the candidate output classes, and this process
enables WE to efficiently employ good classifiers. WE gives significantly
higher performance compared to the DDAG, and requires the number of
classifications on average about half of the number of all possible binary
classifiers.
Additionally, we propose VCF by applying the voting technique to carefully
select the high competitive classes with high confidence. The remaining
candidate classes are recursively eliminated by using WE. Although the number
of classifications of VCF is equal to that of Max Wins, it shows the highest
accuracy compared to all the other algorithms.
Finally, more experiments were conducted to compare our proposed algorithms
and Max Wins in order to find the suitable scenario for using each of them.
The RADAG should be chosen when the number of classes is large and the
classification time is the most concern. VCF shows the highest accuracy among
our proposed algorithms, and it should be selected when the time constraints
is not the main concern. In a general case, WE is the most suitable method
because it is superior to Max Wins in terms of accuracy and time. All of our
techniques apply the generalization performance for organizing the use of the
binary classifiers. This measure can be optimally estimated by the mechanism
of $k$-fold cross-validation that is independent of base learners.
Consequently, all our proposed methods can be also applied to other base
classifiers such as logistic regression, perceptron, linear discriminant
analysis, etc. The estimation of generalization errors using k-fold cross
validation requires additional computation, and this can be thought of as a
drawback of our methods. However, the estimation is done in the offline
training phase, and thus it does not affect the performance in the
classification phase.
## 7 Acknowledgment
The authors would like to thank Dr.Peerapon Vateekul for his valuable comments
on an earlier version of this paper. This research is partially supported by
the Thailand Research Fund, and the Graduate School, Chulalongkorn University.
## References
* [1] V.Vapnik , Statistical Learning Theory, New York, Wiley, 1998.
* [2] V.Vapnik, An overview of statistical learning theory, IEEE Transactions on Neural Networks. 10 (1999) 988-999.
* [3] J.H.Friedman, Another approach to polychotomous classification, Technical report, Stanford University, Department of Statistics, 1996.
* [4] C.C.Chang, L.Chien, and Y.Lee, A novel framework for multi-class classification via ternary smooth support vector machine, Pattern Recognition 44 (2011) 1235-1244.
* [5] J.Manikandan, and B.Venkataramani, Study and evaluation of a multi-class SVM classifier using diminishing learning technique, Neurocomputing 73 (2010) 1676-1685.
* [6] T.G.Dietterich, and G.Bakiri, Solving multiclass learning problems via error-correcting output codes, Journal of Artificial Intelligence Research 2 (1995) 263-286.
* [7] E.L.Allwein, R.E.Schapire, and Y.Singer, Reducing Multiclass to Binary: A Unifying Approach for Margin Classifiers, Journal of Machine Learning Research 1 (2000) 113-141.
* [8] M.A.Bagheri, G.Montazer, and E.Kabir, A subspace approach to error correcting output codes, Pattern Recognition Letters 34 (2013) 176-184.
* [9] L.I.Kuncheva, Using diversity measures for generating error-correcting output codes in classifier ensembles, Pattern Recognition Letters 26 (2005) 83-90.
* [10] A.C.Lorena, and A.C.P.L.F.Carvalho Evaluation functions for the evolutionary design of multiclass Support Vector Macines, International Journal of Computational Intelligence and Applications 8 (2009) 53-68.
* [11] J.Platt, N.Cristianini, and J.Shawe-Taylor, Large margin DAGs for multiclass classification, Proceedings of Neural Information Processing Systems, MIT Press (2000) 547-553.
* [12] C.Hsu, and C.Lin, A comparison of methods for multiclass support vector machines, IEEE Transactions on Neural Networks 13 (2002) 415-425.
* [13] B.Kijsirikul,and N.Ussivakul, Multiclass support vector machines using adaptive directed acyclic graph., Proceedings of International Joint Conference on Neural Networks (IJCNN) (2002) 980-985.
* [14] F.Takahashi, and S.Abe, Optimizing directed acyclic graph support vector machines, Proceedings of Artificial Neural Networks in pattern recognition (2003) 166-170.
* [15] A.C.Lorena, and A.C.P.L.F.Carvalho, Building binary-tree-based multiclass classifiers using separability measures, Neurocomputing 73 (2010) 2837-2845.
* [16] R.Li, A.Li, T.Wang, and L.Li, Vector projection method for unclassifiable region of support vector machine, Expert Systems with Applications 38 (2011) 856-861.
* [17] C.Blake, E.Keogh, and C.Merz, UCI repository of machine learning databases, Department of Information and Computer Science, University of California, Irvine, 1998.
* [18] V.Vapnik, The Nature of Statistical Learning Theory, London, UK, Springer-Verlag, 1995.
* [19] V. N.Vapnik, and A.Y.Chervonenkis, Teoriya Raspoznavaniya Obrazov: Statisticheskie Problemy Obucheniya. (Russian) [Theory of Pattern Recognition: Statistical Problems of Learning]., Moscow: Nauka, 1974.
* [20] P.L.Bartlett, and J.Shawe-Taylor, Generalization performance of support vector machines and other pattern classifiers, Advances in Kernel Methods - Support Vector Learning, MIT Press, Cambridge, USA, (1999) 43-54.
* [21] C.Burges, A tutorial on support vector machines for pattern recognition, Data Mining and Knowledge Discovery 2 (1998) 121-167.
* [22] T.Mitchell, Machine Learning, McGraw Hill, 1997.
* [23] R.Johnson and, G.Bhattacharyya, Statistics: principles and methods, New York, Wiley, 2001.
* [24] W.Cook, and A.Rohe, Computing minimum-weight perfect matchings, Technical Report 97863, Forschungsinstitut f$\ddot{u}$r Diskrete Mathematik, Universit$\ddot{a}$t Bonn, (1999).
* [25] T.Joachims, Making large-scale SVM learning practical, Advances in Kernel Methods - Support Vector Learning, MIT Press (1998).
* [26] T.Joachims, SVMlight, http://ais.gmd.de/~thorsten/svm_light, (1999).
Patoomsiri Songsiri received the B.Sc. degree in Computer Science (First class
honor) from Prince of Songkla University, Thailand, in 2001 and the M.Sc.
degree in Computer Science from Chulalongkorn University, Thailand, in 2006.
She is currently working toward the Ph.D. degree in Computer Engineering at
Chulalongkorn University. Her research interests include Pattern Recognition
and Machine Learning.
Thimaporn Phetkaew received her B.Sc. degree in Applied Mathematics and she
also received her M.Sc. degree in Computer Science from Prince of Songkla
University, Thailand in 1997 and 2000, respectively. In 2004, she received her
Ph.D. degree in Computer Engineering from Chulalongkorn University, Thailand.
Since 2004, as a lecturer, she has been with the School of Informatics,
Walailak University. She is also a member of Informatics Innovation Research
Unit at Walailak University. Her research interests include Data Mining,
Machine Learning, and Software Testing.
Boonserm Kijsirikul received the B.Eng. degree in Electronic and Electrical
Engineering, the M.Sc. degree in Computer Science, and the Ph.D. in Computer
Science from Tokyo Institute of Technology, Japan, in 1986, 1990, and 1993,
respectively. He is currently a Professor at the Department of Computer
Engineering, Chulalongkorn University, Thailand. His current research
interests include Machine Learning, Artificial Intelligence, Natural Language
Processing, and Speech Recognition.
|
arxiv-papers
| 2013-09-11T08:59:07 |
2024-09-04T02:49:50.690866
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Patoomsiri Songsiri, Thimaporn Phetkaew, Boonserm Kijsirikul",
"submitter": "Patoomsiri Songsiri Ms.",
"url": "https://arxiv.org/abs/1309.2765"
}
|
1309.2780
|
# $Z^{0}$-Boson Decays in a Strong Electromagnetic Field 111Published in:
Yad. Fiz. 72 (2009) 1078 , [Phys. At. Nucl. 72 (2009) 1034].
A.V.Kurilin222E-mail address: [email protected]
Moscow Technological Institute 333on leave from Moscow State Open Pedagogical
University,
Leninsky Prospect, 38a, Moscow, 119334, Russia
The probability of $Z^{0}$-boson decay to a pair of charged fermions in a
strong electromagnetic field, $Z^{0}\rightarrow\bar{f}f$, is calculated. On
the basis of a method that employs exact solutions to relativistic wave
equations for charged particles, an analytic expression for the partial decay
width $\Gamma(\varkappa)=\Gamma(Z^{0}\rightarrow\bar{f}f)$ is obtained at an
arbitrary value of the parameter
$\varkappa=eM_{Z}^{-3}\sqrt{-(F_{\mu\nu}q^{\nu})^{2}}$, which characterizes
the external-field strength. The total $Z^{0}$-boson decay width in an intense
electromagnetic field, $\Gamma_{Z}(\varkappa)$, is calculated by summing these
results over all known generations of charged leptons and quarks. It is found
that, in the region of relatively weak fields ($\varkappa<0,06$), the field-
induced corrections to the standard $Z^{0}$-boson decay width in a vacuum do
not exceed 2%. As $\varkappa$ increases, the total decay width
$\Gamma_{Z}(\varkappa)$ develops oscillations against the background of its
gradual decrease to the absolute-minimum point. At $\varkappa_{\rm
min}=0,445$, the total $Z^{0}$-boson decay width reaches the minimum value of
$\Gamma_{Z}(\varkappa_{\rm min})=2,164$ GeV, which is smaller than the
$Z^{0}$-boson decay width in a vacuum by more than 10%. In the region of
superstrong fields ($\varkappa>1$), $\Gamma_{Z}(\varkappa)$ grows
monotonically with increasing external-field strength. In the region
$\varkappa>5$, the $t$-quark-production process $Z^{0}\rightarrow\bar{t}t$,
which is forbidden in the absence of an external field, begins contributing
significantly to the total decay width of the $Z^{0}$\- boson.
## 1\. INTRODUCTION
It is not an exaggeration to say that $Z^{0}$-bosons along with their charged
partners, $W^{\pm}$-bosons, have been at the focus of attention in high-energy
physics over more than the past two decades. The properties of these
particles, which are mediators of weak interaction between leptons and quarks,
were studied in detail in experiments at the LEP and SLC electron-positron
colliders. In a series of experiments performed between 1989 and 1995 at CERN
[1], enormous statistics of observations of $Z^{0}$-boson production and
decays were collected. In all, approximately $2\cdot 10^{6}$ events associated
with leptonic modes of $Z^{0}$-boson decay and more than $1.5\cdot 10^{7}$
events of transformation of these unstable particles into hadrons were
detected. The results of these investigations made it possible to determine,
among other things, the mass of the $Z^{0}$-boson, $M_{Z}=91,1876\pm 0,0021\
\mbox{GeV}$, and its total decay width, $\Gamma_{Z}=2,4952\pm 0,0023\
\mbox{GeV}$, to an extremely high degree of precision [2]. Many other
parameters of the process $e^{+}e^{-}\rightarrow Z^{0}\rightarrow\bar{f}f$
were also analyzed both in the vicinity of the $Z$-resonance and far from it
[3]. The precision of experimental data reached in recent years has given a
motivation to theoretical physicists for performing formidable work on
calculating radiative corrections to the processes involving $Z^{0}$-boson
production and decay (see, for example, [4, 5, 6, 7] and references therein).
Not only does $Z^{0}$-boson physics provide a precision test for Standard
Model predictions concerning particle interactions, but it is also a unique
tool for seeking manifestations of new physics (beyond the Standard Model),
which will replace sooner or later the generally accepted $SU(3)\times
SU(2)\times U(1)$ scheme. Despite the impressive successes of the Standard
Model in adequately describing a formidable set of experimental data, there is
no doubt among the physics community that the modern model of fundamental
interactions cannot be conclusive because of a number of fundamental
theoretical drawbacks. These include, first, an enormous number (more than
19!) of independent input parameters — coupling constants, the masses of quark
and leptons, the parameters of the Cabibbo-Kobayashi-Maskawa mixing matrix,
etc; second, the instability of the mass of the as-yet-undiscovered Higgs
boson with respect to quadratically divergent radiative corrections and the
allied hierarchy problem; and, third, the isolated character of gravitational
interactions, which do not fit in the existing quantum model. This is not the
whole story, however. Paradoxically as it is, the main drawback of the
Standard Model at the present time is that it is in perfect agreement with all
experimental data accumulated thus far, providing no hint as to the nature of
new physics that will have to replace it. This kind of ”recess” in the
experimental development of physics theory caused a flash of creative activity
in theoretical physics in attempts at guessing vague outlines of a new theory
of everything. By no means do Grand Unification theories, supersymmetry,
supergravity, superstrings, supermembranes, and $M$-theory exhaust the list of
ideas that require an experimental verification. The questions that a future
theory should answer include that of the origin of particle masses, that of
why our space-time is four-dimensional, that of why several generations of
leptons exist, that of how many flavors quarks have, and that of how many
neutrino flavors exist. Even today, physicists try to find answers to many of
these questions by analyzing available experimental data.
For example, the total invisible width of $Z^{0}$-bosons, $\Gamma_{\rm
inv}=499,0\pm 1,5\ \mbox{MeV}$, with respect to the decays of these particles
through channels that are rather difficult to detect plays an important role
in $Z^{0}$-boson physics. In the experiments performed at the LEP accelerator,
the invisible width $\Gamma_{\rm inv}$ was determined indirectly as the
difference of the total decay width of the $Z^{0}$-boson and its total width
with respect to all observed decays. In the Standard Model, the invisible
decay width of the $Z^{0}$-boson is interpreted as the partial width with
respect to decays to neutrinos of various flavors:
$\Gamma_{\rm
inv}=\sum_{i}^{N_{\nu}}\Gamma(Z^{0}\rightarrow\nu_{i}\tilde{\nu}_{i})=N_{\nu}\Gamma_{\nu\tilde{\nu}}$
This quantity can be considered as some kind of a counter of neutrino
generations, and its precise measurement is of particular interest for new
physics (see, for example, [8]). It is intriguing that the present-day
experimental value of the number of neutrinos from the invisible decay width
of the $Z^{0}$-boson, $N_{\nu}=2.92\pm 0,06$ [2], differs somewhat from the
result expected in the Standard Model, $N_{\nu}=3$.
Exotic channels of $Z^{0}$-boson decay frequently provide unique tests for
some fundamental principles of modern physics. For example, analysis of the
forbidden (in a vacuum) $Z^{0}$-boson decays to photons
($Z^{0}\rightarrow\gamma\gamma$) and gluons ($Z^{0}\rightarrow gg$) makes it
possible to establish a quantitative measure of the principle of quantum
indistinguishability of integer-spin particles — Bose symmetry for photons and
gluons [9]. According to the Landau-Yang theorem [10], a massive vector
particle cannot decay to massless vector states. In an external
electromagnetic field, however, the reactions $Z^{0}\rightarrow\gamma\gamma$
and $Z^{0}\rightarrow gg$ become possible [11].
The problem of searches for new possible manifestations of supersymmetry is
also tightly related to $Z^{0}$-boson physics. There is still a hope for
finding relatively light superparticles produced in rare channels of
$Z^{0}$-boson decay. As an example, one can consider the reaction
$Z^{0}\rightarrow\tilde{g}\tilde{g}$, in which a pair of light gluinos,
$\tilde{g}$, having a mass of about $m_{\tilde{g}}=12\div 16$ GeV appear. This
process, which is extensively discussed in the literature (see [12] and
references therein), is of special interest in view of prospects for observing
new-physics manifestations at next-generation $p\bar{p}$-colliders. Projects
of experiments with $Z^{0}$-bosons at the Large Hadron Collider(LHC) also
caused interest in studying the properties of these particles in the presence
of a quark-gluon plasma, which is quite an unusual environment of [13].
In view of the aforesaid, it is of particular interest to discuss some other
new phenomena in $Z^{0}$-boson physics that manifest themselves under unusual
conditions. In this study, we investigate changes in partial decay widths of
the $Z^{0}$-boson in the presence of strong electromagnetic fields. Interest
in such investigations is explained both by astrophysical applications and by
applications in the physics of relativistic particles channeling through
single crystals (so-called channeling of particles [14]). It is well known
that the strength of electric fields generated by the axes and planes of
single crystals may reach formidable values of ($E\geq 10^{10}\;$V/m) [15]
over macroscopic distances. At the same time, it is well known that
astrophysics studies exotic objects (neutron stars, white dwarfs) that, at the
latest stages of their evolution, can undergo a strong compression, and this
leads to a significant increase in the magnetic field strength inside them
($H\geq 10^{8}\div 10^{13}\;$ G [16]). At such enormous values of the external
field strength, the physics of quantum processes changes significantly.
External fields frequently remove the forbiddance of specific reactions whose
occurrence is impossible in a vacuum.
The first investigations into the physics of relativistic particles in the
presence of a strong external field were performed within quantum
electrodynamics (see, for example, [17, 18, 19]). Later, the technique of
calculations that was developed there was also used in non-Abelian gauge
theories [20]. Respective calculations rely on the method of exact solutions
to the wave equations for charged particles, which makes it possible to take
into account the interaction with the electromagnetic field beyond standard
perturbation theory. In this case, the wave functions for all charged
particles and their propagators are modified in quite a unwieldy manner,
depending on the configuration of an external field [21]. Although the
resulting expressions are cumbersome and although the calculations are quite
involved, the ultimate results, which are valid at arbitrarily high values of
the external-electromagnetic-field strength, carry information that cannot be
obtained by perturbative methods, and this is an obvious advantage of the
method. From this point of view, investigation of quantum processes in
superstrong fields provides a unique possibility of analyzing the self-
consistency of the physical theory globally in the case where an expansion in
a small parameter is impossible.
The method of calculations that is used below is based on the crossed-field
model, which was successfully applied in the previous studies of the present
author to studying photino-pair production [22] and to analyzing the decays of
$W$-bosons in an external field [23].
## 2\. PROBABILITY FOR $Z^{0}$-BOSON DECAY
IN EXTERNAL FIELDS
In the leading order of perturbation theory in the coupling constants $g_{V}$
and $g_{A}$, the matrix element of $Z^{0}$-boson decay to a pair of charged
leptons, $\ell^{\pm}$, is given by
$S_{fi}=i\int
d^{4}x\overline{\Psi}_{\ell^{-}}(x,p)\gamma^{\mu}(g_{V}+\gamma^{5}g_{A})\Psi_{\ell^{+}}^{c}(x,p^{\prime})Z_{\mu}(x,q).$
(1)
The presence of an external electromagnetic field is taken into account via
choosing specific wave functions $\Psi_{\ell^{-}}(x,p)$ and
$\Psi_{\ell^{+}}(x,p^{\prime})$ for the charged leptons $\ell^{\pm}$. This
approach, which, in particle physics, is referred to as the Furry
representation, yields, for the probability of $Z^{0}$-boson decay, an
expression in which electromagnetic interactions are taken into account
exactly, beyond standard perturbation theory in the electromagnetic coupling
constant $\alpha=e^{2}/4\pi\simeq 1/137$. The explicit form of the wave
functions for charged particles in an external electromagnetic field depends
on the field configuration. In the present study, we restrict ourselves to the
crossed-field configuration, in which case the strength tensor
$F_{\mu\nu}=\partial_{\mu}A_{\nu}-\partial_{\nu}A_{\mu}$ has the form
$F_{\mu\nu}=C(k_{\mu}a_{\nu}-k_{\nu}a_{\mu}).$ (2)
It is assumed in this case that the external-electromagnetic-field potential
$A_{\mu}(x)$ is chosen in the gauge
$A_{\mu}(x)=a_{\mu}C(k_{\nu}x^{\nu}),$ (3)
Here, the unit constant $4$-vector $a_{\mu}$ determines the spatial field
configuration and satisfies the conditions
$a_{\mu}a^{\mu}=-1,\qquad
F_{\mu\nu}=(a_{\mu}F_{\nu\lambda}-a_{\nu}F_{\mu\lambda})a^{\lambda}.$ (4)
A crossed field is a particular case of the field of a plane electromagnetic
wave having an isotropic wave $4$-vector $k_{\mu}\ (k^{2}=0)$. The invariant
parameter $C$ characterizes the crossed-field strength and has the dimensions
of energy. Recall that, in an arbitrary reference frame, a crossed field is a
superposition of constant electric and magnetic fields whose strength vectors
are orthogonal and are equal in absolute value: $E=H=k_{0}C$. Moreover, the
two relativistic invariants of the electromagnetic field are zero:
$F_{\mu\nu}F^{\mu\nu}=F_{\mu\nu}\tilde{F}^{\mu\nu}=0.$ (5)
In the crossed-field model, the choice of wave $4$-vector $k_{\mu}$ is quite
arbitrarily. It is only necessary that this vector be isotropic, have
dimensions of energy, and satisfy the conditions
$k_{\mu}a^{\mu}=0,\qquad F^{\mu\lambda}F_{\lambda\nu}=C^{2}k^{\mu}k_{\nu}.$
(6)
In particular, the wave $4$-vector $k^{\mu}$ can be expressed in terms of the
electromagnetic-field-strength tensor and the constant electromagnetic-
potential $4$-vector $a_{\lambda}$ as
$k^{\mu}=-\frac{1}{C}\cdot F^{\mu\lambda}a_{\lambda}.$ (7)
In dealing with the problem of $Z^{0}$-boson decay, it is advisable to express
the wave vector $k^{\mu}$ in terms of the $Z^{0}$-boson 4-momentum $q^{\nu}$
and the electromagnetic-field-strength tensor $F^{\mu\lambda}$ (see next
section) as
$k^{\mu}=\frac{e^{2}\Delta}{2M_{Z}^{6}\varkappa^{2}}\cdot
F^{\mu\lambda}F_{\lambda\nu}q^{\nu}.$ (8)
In an arbitrary constant uniform electromagnetic field, the probability of the
decay $Z^{0}\rightarrow\ell^{+}\ell^{-}$ is a function of three invariant
dimensionless parameters,
$P(Z^{0}\rightarrow\ell^{+}\ell^{-})=P(\varkappa,a,b)$. These parameters are
determined by the strength tensor of an external macroscopic field as
$\displaystyle\varkappa=\frac{e}{M_{Z}^{3}}\sqrt{-(F_{\mu\nu}q^{\nu})^{2}},$
(9) $\displaystyle a=-\frac{e^{2}}{4M_{Z}^{4}}F_{\mu\nu}\tilde{F}^{\mu\nu},$
(10) $\displaystyle b=\frac{e^{2}}{4M_{Z}^{4}}F_{\mu\nu}F^{\mu\nu}.$ (11)
In the crossed-field model, the parameter $\varkappa$ has the simplest form
$\varkappa=\frac{eC}{M_{Z}^{3}}(q_{\mu}k^{\mu}),$ (12)
the square of the wave vector (8) having the form
$k^{2}=\frac{\Delta^{2}}{4M_{Z}^{2}}\left({a^{2}\over\varkappa^{4}}-{2b\over\varkappa^{2}}\right).$
(13)
In a crossed field, the two parameters (10) and (11) of the external-field
strength vanish by virtue of the equality in (5), $a=b=0$; therefore, the
right-hand side of Eq.(13) also vanishes. For a constant uniform
electromagnetic field of general form, the condition that the wave vector in
(8) is isotropic is not satisfied: $k^{2}\neq 0$. Even in this case, however,
we can use the semiclassical crossed-field model, which ensures a satisfactory
description of the decay probability in an electromagnetic field of general
form, provided that $a\ll\varkappa^{2}$ and $b\ll\varkappa^{2}$. The closer to
zero the right-hand side of Eq.(13) and the higher the precision to which the
condition of isotropicity is satisfied, the more precise the results that are
obtained in this approximation. The above constraints are obviously satisfied
in the region of relatively weak electromagnetic fields for $a\ll 1$ and $b\ll
1$. Thus, the decay probability $P(Z^{0}\rightarrow\ell^{+}\ell^{-})$ in a
crossed field will correspond to the first leading term $P_{0}$ in the
semiclassical expansion of the total decay probability $P(\varkappa,a,b)$ in a
power series in the vectors of the external-field strengths at small values of
the parameters $a$ and $b$; that is,
$P(\varkappa,a,b)=P_{0}(\varkappa,0,0)+a\frac{\partial P}{\partial
a}(\varkappa,0,0)+b\frac{\partial P}{\partial b}(\varkappa,0,0)+\ldots$ (14)
This relation between the probabilities of quantum processes in a crossed
field and a constant electromagnetic field of general form was obtained for
the first time within quantum electrodynamics [18].
The wave functions for charged fermions in the field of a plane
electromagnetic wave were obtained as far back as 1937 by D.M.Volkov (see,for
example, [24]). In the case of a crossed field, these exact solutions to the
Dirac equations for leptons $\ell^{\pm}$ in an external field have the form
$\displaystyle\Psi_{\ell^{-}}(x,p)=\exp\biggl{[}-ipx+\frac{ieC(pa)}{2(pk)}(kx)^{2}-\frac{ie^{2}C^{2}}{6(pk)}(kx)^{3}\biggr{]}\times$
$\displaystyle\times\biggl{\\{}1-\frac{e(kx)}{4(pk)}(F_{\mu\lambda}\gamma^{\mu}\gamma^{\lambda})\biggr{\\}}\frac{u(p)}{\sqrt{2p_{0}V}};$
(15)
$\displaystyle\Psi_{\ell^{+}}^{c}(x,p^{\prime})=\exp\biggl{[}ip^{\prime}x+\frac{ieC(p^{\prime}a)}{2(p^{\prime}k)}(kx)^{2}-\frac{ie^{2}C^{2}}{6(p^{\prime}k)}(kx)^{3}\biggr{]}\times$
$\displaystyle\times\biggl{\\{}1+\frac{e(kx)}{4(p^{\prime}k)}(F_{\mu\lambda}\gamma^{\mu}\gamma^{\lambda})\biggr{\\}}\frac{u^{c}(p^{\prime})}{\sqrt{2p^{\prime}_{0}V}};$
(16)
The spin part of these wave functions normalized to a three-dimensional
spatial volume $V$ is determined by the constant Dirac bispinors $u(p)$ and
$u^{c}(p^{\prime})$. The antilepton wave function
$\Psi_{\ell^{+}}(x,p^{\prime})$ differs from expression (2.) only by the
reversal of the sign of the electric charge $e$: $e\rightarrow-e$. In the
expression for the matrix element of the decay in (1), the antilepton wave
function appears in the charge-conjugate form (2.), which can be derived on
the basis of the Dirac-conjugate bispinor
$\bar{\Psi}_{\ell^{+}}(x,p^{\prime})$ with the aid of the charge-conjugation
matrix $U_{C}=-i\gamma^{0}\gamma^{2}$; that is,
$\Psi^{c}_{\ell^{+}}(x,p^{\prime})=U_{C}\bar{\Psi}^{T}_{\ell^{+}}(x,p^{\prime})$
(17)
As for the $Z^{0}$-boson, its wave function does not change in an
electromagnetic field because it is electrically neutral:
$Z_{\mu}(x,q)=\exp\biggl{(}-iqx\biggr{)}\frac{v_{\mu}(q)}{\sqrt{2q_{0}V}}.$
(18)
The spin states of the $Z^{0}$-boson are characterized by the complex
polarization $4$-vector $v_{\mu}(q)$ satisfying standard conditions for
massive vector fields; that is,
$\displaystyle v_{\mu}(q)q^{\mu}=0,\quad v^{*}_{\mu}(q)v^{\mu}(q)=-1,$ (19)
$\displaystyle\sum_{\sigma=1}^{3}v_{\mu}(q,\sigma)v^{*}_{\nu}(q,\sigma)=-g_{\mu\nu}+q_{\mu}q_{\nu}/M_{Z}^{2}.$
We substitute the particle wave functions (2.), (2.), and (18) into expression
(1) for the $S$-matrix element and perform integration of $\mid
S_{fi}\mid^{2}$ over the phase space of final leptons. Upon summation over the
spin states of the lepton-antilepton pair and averaging over polarizations of
the $Z^{0}$-boson, we arrive at
$\displaystyle
P(Z^{0}\rightarrow\ell^{+}\ell^{-}\mid\varkappa)=\frac{(g_{V}^{2}+g_{A}^{2})M_{Z}^{2}}{12\pi^{2}q_{0}}\int\limits_{0}^{1}du\biggl{\\{}\biggl{[}1-(1+3\lambda)\frac{m_{\ell}^{2}}{M_{Z}^{2}}\biggr{]}\Phi_{1}(z)-$
$\displaystyle-\frac{2\varkappa^{2/3}}{\left[u(1-u)\right]^{1/3}}\biggl{[}1-2u+2u^{2}+(1+\lambda)\frac{m_{\ell}^{2}}{M_{Z}^{2}}\biggr{]}\Phi^{\prime}(z)\biggr{\\}}.$
(20)
The probability of $Z^{0}$-boson decay is expressed in terms of the Airy
functions $\Phi^{\prime}(z)$ and $\Phi_{1}(z)$ (for necessary details
concerning these special mathematical functions, see the Appendix) depending
on the argument
$z=\frac{m_{\ell}^{2}-M_{Z}^{2}u(1-u)}{M_{Z}^{2}[\varkappa u(1-u)]^{2/3}}.$
(21)
The couplings of the $Z^{0}$-boson to the charged leptons of all three
generations, $e^{\pm},\mu^{\pm},$ and $\tau^{\pm}$, can be expressed in terms
of the Weinberg angle $\theta_{\rm W}$ and the electroweak coupling constant
$g=e/\sin\theta_{\rm W}$ as
$g_{V}=-\frac{g(1-4\sin^{2}\theta_{\rm W})}{4\cos\theta_{\rm W}},\qquad
g_{A}=-\frac{g}{4\cos\theta_{\rm W}}.$ (22)
The dimensionless parameter $\lambda$ in expression (2.) is given by
$\lambda=\frac{g_{A}^{2}-g_{V}^{2}}{g_{A}^{2}+g_{V}^{2}}=\frac{1-(1-4\sin^{2}\theta_{\rm
W})^{2}}{1+(1-4\sin^{2}\theta_{\rm W})^{2}},$ (23)
at $\sin^{2}\theta_{\rm W}=0,23$, it is very close to unity ($\lambda\simeq
0,987$).
## 3\. KINEMATICS OF $Z^{0}$-BOSON DECAY
IN A CROSSED FIELD
The fact that the law of energy-momentum conservation for particles involved
in the decay $Z^{0}\rightarrow\ell^{+}\ell^{-}$ does not have a conventional
form,
$q_{\mu}+k_{\mu}=p_{\mu}+p^{\prime}_{\mu}.$ (24)
is a distinctive feature of processes occurring in an external electromagnetic
field. Along with the $Z^{0}$-boson and final-lepton energy-momentum
$4$-vectors $q_{\mu}$, $p_{\mu}(\ell^{-})$, and $p^{\prime}_{\mu}(\ell^{+})$,
expression (24) also involves the wave vector (8) determining the energy
contribution of the external field. It should be noted that, in this case, we
are dealing with asymptotic momenta that charged particles possess at rather
large distances, where the effect of the external field is negligible, rather
than with their dynamical $4$-momenta in a crossed field. This is the meaning
that the parameters $p_{\mu}$ and $p^{\prime}_{\mu}$ appearing in the Volkov
solutions (2.) and (2.) to the Dirac equation in a crossed field have. It is
indeed extremely difficult to observe the kinematics of the decay
$Z^{0}\rightarrow\ell^{+}\ell^{-}$ in the field itself because the charged-
lepton $4$-momenta are not integrals of the motion. As soon as the leptons
escape from the region of the external field, their energies and momenta do
not change any longer, so that one can measure them experimentally. In this
case, the $4$-momenta of all particles satisfy the ordinary kinematical
conditions
$p^{2}=p^{\prime 2}=m_{\ell}^{2},\quad q^{2}=M_{Z}^{2},\quad k^{2}=0,$ (25)
and this makes it possible to represent the isotropic vector $k^{\mu}$ in the
form
$k^{\mu}=\frac{\Delta}{2(q_{\alpha}F^{\alpha\beta}F_{\beta\sigma}q^{\sigma})}\cdot
F^{\mu\lambda}F_{\lambda\nu}q^{\nu},$ (26)
where the parameter $\Delta$ characterizes the energy imbalance of the
reaction $Z^{0}\rightarrow\ell^{+}\ell^{-}$ in the external field,
$\Delta=2(qk)=(p+p^{\prime})^{2}-q^{2}=2m_{\ell}^{2}+2(pp^{\prime})-M_{Z}^{2}$
(27)
Expression (24) arises in a natural way as the argument of the four-
dimensional Dirac delta function in calculating the $S$-matrix element (1). In
performing integration over the phase space of final leptons, there also
occurs summation of the external-field contributions characterized by various
values of the parameter $\Delta$. This is the circumstance that distinguishes
the kinematics of $Z^{0}$-boson decay in an external field from a similar
process involving a real photon, $\gamma+Z^{0}\rightarrow\ell^{+}\ell^{-}$,
where the $4$-momentum $k_{\mu}$ is fixed and is independent of the final-
lepton momenta. The physical meaning of the above formulas can most easily be
understood in the $Z^{0}$-boson rest frame, where the law of energy-momentum
conservation in (24) can be expressed in terms of the following $4$-vectors:
$q^{\mu}=(M_{Z};{\bf 0}),\quad p^{\mu}=(\varepsilon_{1};{\bf p}_{1}),\quad
p^{\prime\mu}=(\varepsilon_{2};{\bf p}_{2}),\quad k^{\mu}=(k^{0};{\bf k}).$
(28)
We chose the axes of the three-dimensional Cartesian coordinate system in such
a way that the vectors of the electric- and magnetic-field strength, ${\bf E}$
and ${\bf H}$, are directed along the $x$ and $y$ axes, respectively. The
direction of the $3$-momentum ${\bf k}$ formed by the spatial components of
the wave $4$-vector $k^{\mu}$ (26) then coincides with the positive direction
of the $z$ axis, which is parallel to the vector product ${\bf[E\times H]}$.
In the Cartesian coordinate system chosen above, the orientation of the lepton
($\ell^{-}$) and antilepton ($\ell^{+}$) $3$-momenta ${\bf p}_{1}$ and ${\bf
p}_{2}$ with respect to the $z$ axis can be described by two azimuthal angles
$\vartheta_{1}$ and $\vartheta_{2}$ ($0\leq\vartheta_{1,2}\leq\pi$). Figure 1
shows schematically the three-dimensional configuration formed by all of the
aforementioned vectors.
Figure 1: Kinematics of the decay $Z^{0}\rightarrow\ell^{+}\ell^{-}$ in a
crossed field in the $Z^{0}$-boson rest frame.
If the decay $Z^{0}\rightarrow\ell^{+}\ell^{-}$ occurs in the vacuum and if
the external electromagnetic field is inoperative, the lepton and antilepton
have momenta that are equal in absolute value but are oppositely directed and
which can be expressed in terms of their masses and the $Z^{0}$-boson mass.
The total energy taken away by the lepton-antilepton pair is shared between
the particles in equal parts and its total amount is equal to the rest energy
of the $Z^{0}$-boson,
$p_{1}=p_{2}=\frac{M_{Z}}{2}\left(1-{4m_{\ell}^{2}\over
M_{Z}^{2}}\right)^{1/2},\quad\varepsilon_{1}=\varepsilon_{2}=M_{Z}/2,\quad\vartheta_{1}+\vartheta_{2}=\pi.$
(29)
In a crossed field, each of the equalities in (29) is in general violated, and
the energy $\varepsilon_{1}$ that is taken away in the decay by the lepton is
not equal to the antilepton energy $\varepsilon_{2}$. Only fulfillment of the
following conservation laws is guaranteed:
$\displaystyle\varepsilon_{1}+\varepsilon_{2}-M_{Z}=p_{1}\cos\vartheta_{1}+p_{2}\cos\vartheta_{2},$
(30) $\displaystyle p_{1}\sin\vartheta_{1}-p_{2}\sin\vartheta_{2}=0,$ (31)
The vectors ${\bf p}_{1}$, ${\bf p}_{2}$, and ${\bf k}$ lie in the same plane,
but, in this case, it is not mandatory that the lepton and the antilepton fly
apart along the same straight line, $\vartheta_{1}+\vartheta_{2}\neq\pi$. The
deficit in the ordinary law of energy-momentum conservation is compensated in
an external field by the isotropic wave vector $k^{\mu}$ (8).
Formula (30) represents the three-dimensional form of the conservation law of
the covariant projections of the $4$-momenta of all particles on to the
direction of the wave vector $(qk)=(pk)+(p^{\prime}k)$, while relation (31) is
the conservation law for the ordinary three-dimensional projections of the
particle momenta on to the axis that is orthogonal to the vector ${\bf k}$ and
which lies in the plane spanned by the $3$-vectors ${\bf p}_{1}$, ${\bf
p}_{2}$, and ${\bf k}$.
We can now interpret physically the variable $u$ appearing in the integrand on
the right-hand side of (2.) and in the argument of the Airy functions in (21).
This variable is associated with the angles $\vartheta_{1}$ and
$\vartheta_{2}$ of divergence of the lepton-antilepton pair with respect to
the $z$ axis in the rest frame of the initial $Z^{0}$-boson:
$u=\frac{(pk)}{(qk)}=\frac{\varepsilon_{1}-p_{1}\cos\vartheta_{1}}{M_{Z}}=1-\frac{\varepsilon_{2}-p_{2}\cos\vartheta_{2}}{M_{Z}}$
(32)
If the decay $Z^{0}\rightarrow\ell^{+}\ell^{-}$ is unaffected by an external
field, then, by virtue of relations (29), the variable $u$ can be represented
in the form
$u={1\over 2}\left[1-\left(1-{4m_{\ell}^{2}\over
M_{Z}^{2}}\right)^{1/2}\cos\vartheta_{1}\right]={1\over
2}\left[1+\left(1-{4m_{\ell}^{2}\over
M_{Z}^{2}}\right)^{1/2}\cos\vartheta_{2}\right].$ (33)
The interval of its values admissible from the point of view of decay
kinematics in a vacuum is determined by the inequality
$u_{1}\leq u\leq u_{2},\quad\mbox{where}\quad u_{1,2}={1\over
2}\left[1\mp\left(1-{4m_{\ell}^{2}\over M_{Z}^{2}}\right)^{1/2}\right]$ (34)
The limiting transition $\varkappa\rightarrow 0$ in (2.) reproduces exactly
well-known results for the ordinary probability of the decay
$Z^{0}\rightarrow\ell^{+}\ell^{-}$ without any field. The kinematically
allowed region of values of the variable $u$ (34) corresponds to negative
values of the argument of the Airy functions in (21); therefore, the limiting
transition to zero values of the external-electromagnetic-field strength is
implemented by means of the formal substitution of the Heaviside step
$\theta$-function for the Airy function $\Phi_{1}(z)$:
$\Phi_{1}(z)\rightarrow\pi\theta(-z)$ (35)
In the limit $\varkappa\rightarrow 0$, the differential probability of
$Z^{0}$-boson decay with respect to the variable $u$ does not depend on the
direction of divergence of leptons and has the form
$\frac{dP}{du}(Z^{0}\rightarrow\ell^{-}\ell^{+})=\frac{M_{Z}^{2}}{12\pi
q_{0}}\left[g_{V}^{2}+g_{A}^{2}+{2m_{\ell}^{2}\over
M_{Z}^{2}}(g_{V}^{2}-2g_{A}^{2})\right]$ (36)
As a matter of fact, integration with respect to the invariant variable $u$
now reduces to multiplying expression (36) by the phase space of this
variable, $(u_{2}-u_{1})$. As a result, we obtain the well-known formula for
the leptonic-decay width of the $Z^{0}$-boson in a vacuum,
$\Gamma(Z^{0}\rightarrow\ell^{+}\ell^{-})=\frac{G_{\rm
F}M_{Z}^{3}}{12\pi\sqrt{2}}\sqrt{1-4\delta_{\ell}^{2}}\biggl{(}1-c_{\ell}-\delta_{\ell}^{2}(1+2c_{\ell})\biggr{)},$
(37)
where we have introduced the notation
$\delta_{\ell}={m_{\ell}\over M_{Z}},\qquad c_{\ell}=4\sin^{2}\theta_{\rm
W}-8\sin^{4}\theta_{\rm W}\approx 0,497$ (38)
In a crossed electromagnetic field, the kinematics of $Z^{0}$-boson decay
changes, first of all, owing to the increase in the phase space of possible
states of final leptons $\ell^{\pm}$. In the three-dimensional space of the
momenta $(p_{x},p_{y},p_{z})$ of one of the leptons — for example, $\ell^{-}$
— the kinematically allowed region of the decay
$Z^{0}\rightarrow\ell^{+}\ell^{-}$ in a vacuum [see (29)] can be represented
as a three-dimensional sphere, each point of this sphere corresponding to one
of the final states of this lepton. In an external electromagnetic field, all
possible final states of the lepton $\ell^{-}$ already fill some volume in
momentum space. At a fixed value of the azimuthal angle $\vartheta_{1}$, the
momentum of the lepton $\ell^{-}$ can take any values from the interval
$\displaystyle 0\leq p_{1}<p_{\rm max}(\vartheta_{1}),\quad\mbox{where}$
$\displaystyle\ p_{\rm
max}(\vartheta)=\frac{M_{Z}}{\sin^{2}\vartheta}\left(\cos\vartheta+\sqrt{1-\delta_{\ell}^{2}\sin^{2}\vartheta}\right).$
(39)
The energy $\varepsilon_{1}$ of this particle will be bounded by the
inequality
$\displaystyle m_{\ell}\leq\varepsilon_{1}<\varepsilon_{\rm
max}(\vartheta_{1}),\quad\mbox{where}$ $\displaystyle\varepsilon_{\rm
max}(\vartheta)=\frac{M_{Z}}{\sin^{2}\vartheta}\left(1+\cos\vartheta\sqrt{1-\delta_{\ell}^{2}\sin^{2}\vartheta}\right).$
(40)
The variables $\varepsilon_{1}$ and $\vartheta_{1}$ determine unambiguously
decay kinematics in a crossed field. The energy $\varepsilon_{2}$ and the
azimuthal angle $\vartheta_{2}$ of the antilepton $\ell^{+}$ can be expressed
in terms of these variables with the aid of the kinematical relations (30) and
(31) and the standard relativistic relations
$\varepsilon_{1}=\sqrt{p_{1}^{2}+m_{\ell}^{2}},\quad\varepsilon_{2}=\sqrt{p_{2}^{2}+m_{\ell}^{2}}.$
(41)
After some simple algebra, we obtain
$\displaystyle\varepsilon_{2}=M_{Z}-\varepsilon_{1}+\frac{M_{Z}(\varepsilon_{1}-M_{Z}/2)}{M_{Z}-\varepsilon_{1}+p_{1}\cos\vartheta_{1}}$
(42)
$\displaystyle\vartheta_{2}=\arccos\left(\frac{\varepsilon_{1}+\varepsilon_{2}-M_{Z}-p_{1}\cos\vartheta_{1}}{\sqrt{\varepsilon_{2}^{2}-m_{\ell}^{2}}}\right)$
(43)
An analysis of the resulting formulas shows that, as a rule, the leptons
$\ell^{\pm}$ diverge at arbitrary angles $\vartheta_{1}$ and $\vartheta_{2}$
with respect to the $z$ axis. This is true for almost all of the azimuthal
angles, with the exception of the case where the lepton and antilepton move in
the direction parallel to the $z$ axis. In the case of $\sin\vartheta_{1}\neq
0$ and $\sin\vartheta_{2}\neq 0$, the lepton and antilepton energies
$\varepsilon_{1}$ and $\varepsilon_{2}$, respectively, can be expressed in
terms of the emission angles $\vartheta_{1}$ and $\vartheta_{2}$ as
$\displaystyle\varepsilon_{1}\approx
p_{1}\simeq\frac{M_{Z}\sin\vartheta_{2}}{\sin\vartheta_{1}+\sin\vartheta_{2}-\sin(\vartheta_{1}+\vartheta_{2})}$
(44) $\displaystyle\varepsilon_{2}\approx
p_{2}\simeq\frac{M_{Z}\sin\vartheta_{1}}{\sin\vartheta_{1}+\sin\vartheta_{2}-\sin(\vartheta_{1}+\vartheta_{2})}$
(45)
Figure 2: Lepton energy $\varepsilon_{1}$ as a function of the antilepton
emission angle $\vartheta_{2}$ with respect to the $z$ axis at
$\vartheta_{1}=2\pi/3$.
These relations were obtained in the relativistic approximation, where
$\varepsilon_{1}\gg m_{\ell}$, $\varepsilon_{2}\gg m_{\ell}$ and where the
parameter $\delta_{\ell}$ (38) can be disregarded. We can now draw qualitative
conclusions on a typical behavior of the energies of the lepton-antilepton
pair almost over the entire azimuthal-angle range $0<\vartheta_{1,2}<\pi$. It
was indicated above that, at a fixed value of the emission angle
($\vartheta_{1}={\rm const}$) of the lepton $\ell^{-}$, its energy can vary
from the rest energy to the maximum value $\varepsilon_{\rm
max}(\vartheta_{1})$ (3.). Concurrently, the emission angle $\vartheta_{2}$ of
the antilepton $\ell^{+}$ can take an arbitrary value between zero and $\pi$.
As for its energy $\varepsilon_{2}$, one can see that, as the energy
$\varepsilon_{1}$ increases, $\varepsilon_{2}$ first decreases from the value
$\varepsilon_{0}=\frac{M_{Z}}{2}\frac{(1-2\delta_{\ell}+2\delta_{\ell}^{2})}{(1-\delta_{\ell})}\approx\frac{M_{Z}}{2}\left(1-\delta_{\ell}\right)$
(46)
to its local minimum
$\varepsilon_{2}=\frac{M_{Z}\cos\left(\vartheta_{1}/2\right)}{1+\cos\left(\vartheta_{1}/2\right)},$
(47)
whereupon it increases monotonically to infinity as the angle $\vartheta_{2}$
tends to zero, while the lepton energy $\varepsilon_{1}$ tends to its maximum
value $\varepsilon_{\rm max}(\vartheta_{1})$. The above relationship between
the energies of the lepton-antilepton pair and the directions of lepton and
antilepton emission with respect to the $z$ axis is illustrated by the graphs
in Figs. 2 and 3. It should be noted that the minimum value of the antilepton
energy $\varepsilon_{2}$ (47) is reached at the emission angle of
$\vartheta_{2}=\pi-\vartheta_{1}/2$, in which case the energy of the other
particle, $\ell^{-}$,is
$\varepsilon_{1}=\frac{M_{Z}/2}{1+\cos\left(\vartheta_{1}/2\right)}.$ (48)
Figure 3: Antilepton energy $\varepsilon_{2}$ versus the antilepton emission
angle $\vartheta_{2}$ with respect to the $z$ axis at a fixed value of the
lepton emission angle ($\vartheta_{1}=2\pi/3$).
If the lepton escapes in the direction parallel to the $z$ axis, expressions
(44) and (45) are inapplicable, so that one must use the exact kinematical
relations (30) and (31) or their consequences (42) and (43). In this case, two
situations are possible — that in which the leptons fly apart in opposite
directions and that in which the leptons move in the same direction.
In the first case (at $\vartheta_{1}=0$ and $\vartheta_{2}=\pi$ or at
$\vartheta_{1}=\pi$ and $\vartheta_{2}=0$), the energy of the lepton moving
along the positive direction of the $z$ axis can take any value from the rest
energy $\varepsilon_{\rm min}=m_{\ell}$ to infinity. The energy of the second
lepton, which flies in the opposite direction, is virtually constant — it
ranges between the initial value $\varepsilon_{0}$ (46) and its maximum value
$\varepsilon_{\pi}=\lim_{\vartheta\rightarrow\pi}\varepsilon_{\rm
max}(\vartheta)=\frac{M_{Z}}{2}\left(1+\delta_{\ell}^{2}\right).$ (49)
which is very close to $\varepsilon_{0}$. The case where both particles
$\ell^{+}$ and $\ell^{-}$ produced in $Z^{0}$-boson decay move in the same
direction opposite to the positive direction of the $z$ axis
($\vartheta_{1}=\vartheta_{2}=\pi$) is also possible. In this case, their
energies $\varepsilon_{1}$ and $\varepsilon_{2}$ can vary from the rest energy
to the limiting value $\varepsilon_{0}$ (46). In addition, the following
relation can be obtained for this case from (30) in the relativistic
approximation (for $\delta_{\ell}\rightarrow 0$):
$\varepsilon_{1}+\varepsilon_{2}=\frac{M_{Z}}{2}.$ (50)
Thus, one can see that, in relation to the kinematics of the process
$Z^{0}\rightarrow\ell^{+}\ell^{-}$ in a vacuum [see(29)], the physics of
$Z^{0}$-boson decay in an external electromagnetic field is richer in the
number of possible final lepton states. From the mathematical point of view,
this manifests itself in the fact that all values of the invariant variable
$u$ (32) from the domain $0<u<1$ are now admissible. Moreover, it is
noteworthy that, in a crossed field, the physical meaning of this variable
changes because it now characterizes two independent kinematical quantities
simultaneously - for example, $\varepsilon_{1}$ and $\vartheta_{1}$.
Disregarding the cases in which leptons fly apart in the direction parallel to
the $z$ axis, we can express the variable $u$ in terms of the azimuthal angles
$\vartheta_{1}$ and $\vartheta_{2}$ alone by using formulas (44) and (45). The
result is
$u\simeq\frac{\tan({\vartheta_{1}/2})}{\tan({\vartheta_{1}/2})+\tan({\vartheta_{2}/2})}.$
(51)
Figure 4: Differential width with respect to the decay
$Z^{0}\rightarrow\ell^{-}\ell^{+}$ as a function of the invariant angular
variable $u$ (32) in a weak crossed field (at $\varkappa=0,01$)
An analysis of the integrand on the right-hand side of (2.) shows that, in an
external field, not all of the directions of divergence of the lepton and
antilepton are equiprobable. The graphs in Fig. 4 that represent the
dependence of the differential decay width of the $Z^{0}$-boson on the
invariant variable $u$ corroborate that, even in a relatively weak external
field, there are noticeable deviations from a uniformly distributed
probability of the leptonic decay mode in a vacuum [see (36)]. In the case
where leptons produced in the process $Z^{0}\rightarrow\ell^{+}\ell^{-}$ fly
apart at specific angles $(\vartheta_{1,2}\neq 0,\pi)$ in diametrically
opposite directions $(\vartheta_{1}+\vartheta_{2}\approx\pi)$, the invariant
variable $u$ takes values of $u\approx 0,5$. In this region, which is
kinematically allowed for $Z^{0}$-boson decay in a vacuum, the deviations of
the differential decay width $\Gamma^{\prime}_{u}(u,\varkappa)$ from the
corresponding static value of $\Gamma^{\prime}_{u}(u,0)\approx 83$ MeV are in
fact unobservable. If one of the leptons flies away at a specific angle
$\vartheta$ with respect to the $z$ axis and at an energy close to the maximum
value $\varepsilon_{\rm max}(\vartheta)$ in (3.) and if the second lepton
moves along the direction nearly parallel to this axis (in the positive
direction), the variable $u$ takes values in the region around $u\approx 0$ or
$u\approx 1$. In this case, one can observe oscillations of the differential
decay width of the $Z^{0}$-boson that are characterized by a growing
amplitude, which increases and decreases the static vacuum value in (36) by
more than 15%. Similar oscillations of probabilities of quantum processes in
an external electromagnetic field are characteristic of all particle
transformations allowed in a vacuum and were described in detail in the
literature [18, 19].
## 4\. LEPTONIC DECAY WIDTH OF THE $Z^{0}$-BOSON
The integral representation of the partial width of the $Z^{0}$-boson with
respect to its decay to a pair of charged leptons in an external
electromagnetic field can be derived from expression (2.) for the probability
by going over to the rest frame ($q_{0}=M_{Z}$) and by making the change of
the integration variable $x=u(1-u)$. The result is
$\displaystyle\Gamma(Z^{0}\rightarrow\ell^{+}\ell^{-}\mid\varkappa)$
$\displaystyle=$
$\displaystyle\frac{(g_{V}^{2}+g_{A}^{2})M_{Z}}{6\pi^{2}}\int\limits_{0}^{1/4}\frac{dx}{\sqrt{1-4x}}\biggl{\\{}\biggl{[}1-\delta_{\ell}^{2}(1+3\lambda)\biggr{]}\Phi_{1}(z)-$
(52) $\displaystyle-$ $\displaystyle
2\left(\frac{\varkappa^{2}}{x}\right)^{1/3}\biggl{[}1-2x+\delta_{\ell}^{2}(1+\lambda)\biggr{]}\Phi^{\prime}(z)\biggr{\\}},$
where the argument of the Airy functions is now determined by the expression
$z=\frac{\delta_{\ell}^{2}-x}{(\varkappa x)^{2/3}}.$ (53)
Let us now consider the asymptotic estimates obtained from this formula at
various values of the external-field-strength parameter $\varkappa$ (12).
In weak electromagnetic fields such that $\varkappa\ll\delta_{\ell}^{2}\ll 1$,
the partial decay width of the $Z^{0}$-boson can be written as the sum of two
terms,
$\Gamma(Z^{0}\rightarrow\ell^{+}\ell^{-}\mid\varkappa)=\Gamma(Z^{0}\rightarrow\ell^{+}\ell^{-})+\Delta\Gamma(\varkappa).$
(54)
The first term in this expression is the $Z^{0}$-boson width with respect to
the leptonic decay $Z^{0}\rightarrow\ell^{+}\ell^{-}$ in a vacuum [see
expression (37) above]. The experimental value of this quantity is known to a
precision of 0.1%: $\Gamma(Z^{0}\rightarrow\ell^{+}\ell^{-})=83,984\pm 0,086\
\mbox{MeV}$ [2]. The second term characterizes the effect of the external
electromagnetic field. In the leading order of the expansion in the small
parameters $\varkappa$ and $\delta_{\ell}$, this term can be represented in
the form
$\displaystyle\Delta\Gamma(\varkappa)=-\frac{G_{\rm
F}M_{Z}^{3}}{4\pi\sqrt{6}}\Biggl{[}\
\delta_{\ell}^{2}(1-6c_{\ell})\varkappa\cos\left({1\over 3\varkappa}\right)-$
$\displaystyle-(1-c_{\ell})\varkappa^{2}\sin\left({1\over
3\varkappa}\right)+\frac{8\varkappa^{2}}{\sqrt{3}}(1-c_{\ell})\ \Biggr{]}$
(55)
Similar calculations of the probability of $Z^{0}$-boson decay in a magnetic
field were performed in [25], where the contribution of the virtual electron-
positron loop to the amplitude for elastic $Z^{0}$-boson scattering was
considered, its imaginary part being directly related to the
$Z^{0}\rightarrow\ell^{+}\ell^{-}$ decay width being studied. It is worth
noting, however, that the results of the present calculations [expressions
(52) and (4.)] differ somewhat from their counter parts in [25].
For the parameter of the external-electromagnetic-field strength, we will now
consider the region specified by the condition
$\varkappa\gg\delta_{\ell}^{2}$. We use the fact that the masses of the
leptons of all three generations are much smaller than the $Z^{0}$-boson mass
and go over to the limit $\delta_{\ell}\rightarrow 0$ in formulas (52) and
(53). We can then calculate analytically the integral of the Airy function on
the right-hand side of (52) at arbitrary values of the external-field-strength
parameter $\varkappa$. The final expression for the leptonic-decay width of
the $Z^{0}$-boson is written in terms of the Bessel functions $J_{\nu}(w)$
carrying the noninteger indices $\nu=\pm 1/6$ and $\nu=\pm 5/6$ and depending
on the argument $w=(6\varkappa)^{-1}$. We now consider a universal function
$R(\varkappa)$ that describes the degree of influence of the external field on
the leptonic mode of $Z^{0}$-boson decay:
$R(\varkappa)=\frac{\Gamma(Z^{0}\rightarrow\ell^{+}\ell^{-}\mid\varkappa)}{\Gamma(Z^{0}\rightarrow\ell^{+}\ell^{-})}.$
(56)
Figure 5: Oscillations of the partial width of the $Z$-boson with respect to
its decay to a pair of charged leptons in a weak electromagnetic field.
In the massless-lepton approximation ($\varkappa\gg\delta_{\ell}^{2}$), this
function, which is a relative width of the $Z^{0}$-boson with respect to the
decay $Z^{0}\rightarrow\ell^{+}\ell^{-}$ in an electromagnetic field, can be
represented in the form
$R(\varkappa)=\frac{1}{3}\Biggl{[}1+R_{1}(\varkappa)+R_{2}(\varkappa)+R_{3}(\varkappa)\Biggr{]},$
(57)
where the following notation has been introduced:
$R_{1}(\varkappa)=\frac{\pi(1+3\varkappa^{2})}{12\varkappa}\left[J^{2}_{1/6}\left({1\over
6\varkappa}\right)+J^{2}_{-1/6}\left({1\over 6\varkappa}\right)\right],$ (58)
$R_{2}(\varkappa)=\frac{\pi(1-15\varkappa^{2})}{12\varkappa}\left[J^{2}_{5/6}\left({1\over
6\varkappa}\right)+J^{2}_{-5/6}\left({1\over 6\varkappa}\right)\right],$ (59)
$\displaystyle R_{3}(\varkappa)$ $\displaystyle=$
$\displaystyle\frac{\pi(2+15\varkappa^{2})}{6}\Biggl{[}\ J_{1/6}\left({1\over
6\varkappa}\right)J_{-5/6}\left({1\over 6\varkappa}\right)-$ (60)
$\displaystyle-$ $\displaystyle J_{-1/6}\left({1\over
6\varkappa}\right)J_{5/6}\left({1\over 6\varkappa}\right)\ \Biggr{]}.$
Figure 6: Relative $Z^{0}\rightarrow\ell^{+}\ell^{-}$ decay width
$R(\varkappa)$ (56) as a function of the parameter $\varkappa$ at the boundary
of the region of oscillations.
The graphs representing the dependence of the relative leptonic-decay width
(56) normalized to unity at zero external field on the invariant field-
strength parameter $\varkappa$ (9) are on display in Figs. 5, 6, 7. In
relatively weak fields (such that $\varkappa<0,06$), there arise oscillations
of the probability of the decay $Z^{0}\rightarrow\ell^{+}\ell^{-}$, their
amplitude increasing quadratically with the parameter $\varkappa$. In this
region, one can use the asymptotic expansion for the Bessel functions and
approximate formula (57) by the expression
Figure 7: Leptonic-decay width of the $Z^{0}$-boson as a function of the
external-field-strength parameter $\varkappa$ (12) in the vicinity of the
absolute-minimum point.
$\displaystyle
R(\varkappa)=1-\frac{16}{3}\varkappa^{2}+\frac{3\varkappa^{2}}{\sqrt{3}}\sin\left({1\over
3\varkappa}\right)+\frac{65\varkappa^{3}}{\sqrt{3}}\cos\left({1\over
3\varkappa}\right)+$
$\displaystyle+\frac{880\varkappa^{4}}{3\sqrt{3}}\sin\left({1\over
3\varkappa}\right)-\frac{160}{3}\varkappa^{4}+\ldots$ (61)
It is obvious from Figs. 5 and 6 that the maximum deviation of the decay width
oscillating in an external field from its vacuum value in (37) does not exceed
2%. A further increase in the parameter of the external-electromagnetic-field
strength to values in the region $\varkappa>0,06$ leads to a complete
disappearance of oscillations and a monotonic decrease in the relative width
$R(\varkappa)$ down to the absolute minimum $R_{\min}=0,843$ in the vicinity
of the point $\varkappa_{\rm min}=0,213$ (see Fig. 7). In the region of strong
fields (such that $\varkappa>0,213$), the relative width of the $Z^{0}$-boson
with respect to its leptonic decay mode grows monotonically with the field
strength. In the region $\varkappa\gg 1$, this growth can be approximated by
the asymptotic expression.
$R(\varkappa)=\frac{15\
\Gamma^{4}(2/3)}{14\pi^{2}}(3\varkappa)^{2/3}+\frac{1}{3}+\frac{3\
\Gamma^{4}(1/3)}{110\pi^{2}\ (3\varkappa)^{2/3}}+\ldots$ (62)
## 5\. HADRONIC MODES OF $Z^{0}$-BOSON DECAY
The results obtained in the preceding section admit a straight forward
generalization to the case of the decay $Z^{0}\rightarrow\bar{f}f$, where, by
$f$, we denote an arbitrarily charged fermion, $q_{f}$ standing for the
absolute value of its electric charge. As is well known, the charge is
$q_{f}=2/3$ for the up quarks ($u,c,t$) and $q_{f}=1/3$ for the down quarks
($d,s,b$). It is necessary to consider that, for the $Z\bar{f}f$ vertex, the
coupling constants $g_{V}$ and $g_{A}$ (1) also depend on the absolute value
of the electric charge ($q_{f}$). In the Glashow-Weinberg-Salam model of
electroweak interactions, they are determined by the relations
$\mid g_{V}(f)\mid=\frac{g(1-4q_{f}\sin^{2}\theta_{\rm W})}{4\cos\theta_{\rm
W}},\qquad\mid g_{A}(f)\mid=\frac{g}{4\cos\theta_{\rm W}}.$ (63)
Upon making the formal substitution $e\rightarrow q_{f}e$ or
$\varkappa\rightarrow q_{f}\varkappa$ in (2.), the differential
$Z^{0}\rightarrow\bar{f}f$ decay width as a function of the invariant
variables $\varkappa$ (9) and $u$ (32) assumes the form
$\displaystyle\frac{d\Gamma}{du}(Z^{0}\rightarrow\bar{f}f)=\frac{G_{\rm
F}M_{Z}^{3}}{12\pi^{2}\sqrt{2}}\Biggl{\\{}\left[1-c_{f}-\delta_{f}^{2}(1+2c_{f})\right]\Phi_{1}(z_{f})-$
(64)
$\displaystyle-2\left[\frac{q_{f}^{2}\varkappa^{2}}{u(1-u)}\right]^{1/3}\left[(1-c_{f})(1-2u+2u^{2})+\delta_{f}^{2}\right]\Phi^{\prime}(z_{f})\Biggr{\\}}.$
In this expression, use is made of the same notation as in (21) and (38):
$\delta_{f}={m_{f}\over M_{Z}},\qquad c_{f}=4q_{f}\sin^{2}\theta_{\rm
W}(1-2q_{f}\sin^{2}\theta_{\rm W}),$ (65)
$z_{f}=\frac{\delta_{f}^{2}-u(1-u)}{[q_{f}\varkappa u(1-u)]^{2/3}}.$ (66)
The limiting transition to an infinitely weak external field
($\varkappa\rightarrow 0$) leads to reproducing the well-known results for the
partial widths of the $Z^{0}$-boson with respect to its decay to light quarks;
that is,
$\Gamma(Z^{0}\rightarrow\bar{u}u)=\frac{G_{\rm
F}M_{Z}^{3}}{4\pi\sqrt{2}}\left[1-\frac{8}{3}\sin^{2}\theta_{\rm
W}+\frac{32}{9}\sin^{4}\theta_{\rm W}\right].$ (67)
$\Gamma(Z^{0}\rightarrow\bar{d}d)=\frac{G_{\rm
F}M_{Z}^{3}}{4\pi\sqrt{2}}\left[1-\frac{4}{3}\sin^{2}\theta_{\rm
W}+\frac{8}{9}\sin^{4}\theta_{\rm W}\right].$ (68)
These relations are obtained from (37) upon the formal substitution
$c_{\ell}\rightarrow c_{f}$ and $\delta_{\ell}\rightarrow\delta_{f}$ and the
multiplication by the common factor $N_{c}=3$ (number of color quark states).
The inclusion of radiative corrections leads to various modifications of
expressions (67) and (68), whose explicit form depends on the choice of input
parameters in electroweak theory and renormalization scheme for divergent loop
diagrams [4, 7]. For example, Sirlin, who was the author of one of the
pioneering studies in these realms [26], proposed the parametrization
$G_{\rm F}=\frac{\pi\alpha}{\sqrt{2}M_{W}^{2}\sin^{2}\theta_{\rm
W}}\frac{1}{(1-\Delta r)},\qquad\sin^{2}\theta_{\rm
W}=1-\frac{M_{W}^{2}}{M_{Z}^{2}},$ (69)
which has been widely used since then. Here, $\Delta r$ includes all radiative
corrections. For the sake of simplicity, we restrict ourselves to the case of
massless quarks ($\delta_{f}\ll\varkappa$) and, without going in to details of
the calculation of radiative corrections, try to analyze the effect of
external electromagnetic fields on the hadronic decay modes of the
$Z^{0}$-boson. Without allowance for the external-field effect, the width
$\Gamma_{had}=\sum^{5}_{i=1}\Gamma(Z^{0}\rightarrow\bar{q}_{i}q_{i})\simeq
2\Gamma(Z^{0}\rightarrow\bar{u}u)+3\Gamma(Z^{0}\rightarrow\bar{d}d)$
agrees with the experimentally measured value of $\Gamma(Z^{0}\rightarrow
hadrons)=1744,4\pm 2,0$ MeV [2] to a high degree of precision. We use this
approximation in the region of relatively strong electromagnetic fields
($\varkappa\sim 1$) and estimate the width of the $Z^{0}$-boson with respect
to the hadronic decay mode as
$\Gamma_{had}(\varkappa)=2\Gamma(Z^{0}\rightarrow\bar{u}u)R\left(\frac{2}{3}\varkappa\right)+3\Gamma(Z^{0}\rightarrow\bar{d}d)R\left(\frac{1}{3}\varkappa\right),$
(70)
Figure 8: Hadronic-decay width of the $Z^{0}$-boson in an external
electromagnetic field as a function of the parameter $\varkappa$ (12).
where the function $R(\varkappa)$ is determined by expressions (56) – (60).
Figure 8 shows the results of numerical calculations performed for the total
width of the $Z^{0}$-boson with respect to hadronic decay modes in an external
field on the basis of expression (70) with allowance for the production of
only five quarks ($u,d,s,c,$ and $b$). In this region, the contribution of the
heavy $t$ quark is insignificant. This contribution will be considered
individually in the next section.
One can see from Fig. 8 that, in just the same way as the width with respect
to leptonic decay modes (see Fig. 7), the width with respect to the hadronic
modes of $Z^{0}$-boson decay in an intense external field behaves
nonmonotonically. As the parameter $\varkappa$ increases, the width
$\Gamma_{had}(\varkappa)$ first decreases, reaching the absolute minimum of
$\Gamma_{had}(\varkappa_{\rm min})=1497$ MeV at the point $\varkappa_{\rm
min}=0,501$, where upon it begins growing. In a weak electromagnetic field
($\varkappa\ll 1$), the gradual decrease in the hadronic-decay width
$\Gamma_{had}(\varkappa)$ is accompanied by quite complicated oscillations.
## 6\. TOTAL DECAY WIDTH OF THE $Z^{0}$-BOSON AND $t$-QUARK CONTRIBUTION
In a vacuum, $Z^{0}$-boson decay to $t$ quarks via the process
$Z^{0}\rightarrow\bar{t}t$ is forbidden by the energy-conservation law because
the mass of the heavy $t$ quark is rather high, $m_{t}=174,2\pm 3,3$ GeV [2].
The presence of an external electromagnetic field removes this forbiddance,
opening new channels of $Z^{0}$-boson decay. In principle, the production of
arbitrary charged fermions whose masses satisfy the relation
$\delta_{f}=m_{f}/M_{Z}>1/2$ becomes possible. The missing energy for these
reactions comes from the external field, and this resembles tunneling through
a potential barrier in nonrelativistic quantum mechanics. At small values of
the electromagnetic-field strength (such that $\varkappa\ll 1$), the
$Z^{0}\rightarrow\bar{f}f$ decay width is exponentially small:
Figure 9: Partial width of the $Z^{0}$-boson with respect to its decay to
$t$-quarks in an external electromagnetic field (for the description of the
curves, see the main body of the text).
$\Gamma(Z^{0}\rightarrow\bar{f}f)=\frac{G_{\rm
F}M_{Z}^{3}\left(8\delta_{f}^{2}+1-6c_{f}\right)\delta_{f}^{2}q_{f}}{8\pi\sqrt{6}\left(4\delta_{f}^{2}-1\right)\sqrt{8\delta_{f}^{2}+1}}\varkappa\
\exp\left[-\frac{\left(4\delta_{f}^{2}-1\right)^{3/2}}{3q_{f}\varkappa}\right].$
(71)
This formula was derived from expression (64) upon integration with respect to
the invariant variable $u$ (32) in the approximation $z_{f}\gg 1$ (66).
Substituting the $t$-quark charge $q_{f}=2/3$ and performing summation over
three color states, we obtain an asymptotic estimate for the partial width of
the $Z^{0}$-boson with respect to its decay to $t$-quarks in a relatively weak
electromagnetic field; that is,
$\displaystyle\Gamma(Z^{0}\rightarrow\bar{t}t)=\frac{G_{\rm
F}m_{t}^{2}M_{Z}^{2}\left[8m_{t}^{2}+\left(1-16\sin^{2}\theta_{\rm
W}+\frac{64}{3}\sin^{4}\theta_{\rm
W}\right)M_{Z}^{2}\right]}{4\pi\sqrt{6}\left(4m_{t}^{2}-M_{Z}^{2}\right)\sqrt{8m_{t}^{2}+M_{Z}^{2}}}\times$
$\displaystyle\varkappa\
\exp\left[-\frac{\left(4m_{t}^{2}-M_{Z}^{2}\right)^{3/2}}{2\varkappa
M_{Z}^{3}}\right].$ (72)
The results of numerical calculations show that the above approximate formula
for $\Gamma(Z^{0}\rightarrow\bar{t}t)$ leads to an error not greater than 4%
for all values of the external-field-strength parameter from the region
$\varkappa\leq 5$. As the parameter $\varkappa$ grows, the deviation of the
results obtained according to the asymptotic formula (6.) from the precise
value of the partial width of the $Z^{0}$-boson with respect to the decay
$Z^{0}\rightarrow\bar{t}t$ in an external field increases [see Fig.9, where
the dotted curve corresponds to the calculation according to formula (6.),
while the solid curve represents the results of a numerical integration of
expression (64)].
Figure 10: Total decay width of the $Z^{0}$-boson in an external
electromagnetic field (for the description of the curves, see the main body of
the text).
The $t$-quark contribution to the total decay width of the $Z^{0}$-boson
remains negligible (less than 1%) up to strength-parameter values of about
$\varkappa=6,4$. Figure 10 displays the total decay width of the $Z^{0}$-boson
as a function of the parameter $\varkappa$ in the region of strong fields
($\varkappa\leq 10$). The solid curve corresponds to the total decay width of
the $Z^{0}$-boson with allowance for the $t$-quark contribution, while the
dotted curve represents the same quantity calculated without taking into
account the $t$-quark contribution. One can clearly see a trend toward the
growth of both the absolute value of the total decay width of the
$Z^{0}$-boson and the fraction of the $t$-quark contribution in this decay
width. As the external-field strength increases, the process
$Z^{0}\rightarrow\bar{t}t$, which is forbidden in a vacuum, becomes dominant
in superstrong electromagnetic fields. By way of example, it can be indicated
that, in the $Z^{0}$-boson decay width at $\varkappa=100$, the fraction
associated with $t$-quarks is as large as 50%. This circumstance clearly
demonstrates how an external electromagnetic field can change drastically the
standard physics of quantum processes in a vacuum and serve as some kind of a
catalyst for a number of new and nontrivial phenomena.
## 7\. CONCLUSIONS
The effect of a strong electromagnetic field on the probability of
$Z^{0}$-boson decay has been investigated. The present calculations have been
performed within the crossed-field model, which makes it possible to trace
characteristic variations in the modes of $Z^{0}$-boson decay to known leptons
and quarks. We have found that, in the region of relatively strong fields
($\varkappa\sim 1$), all partial decay widths
$\Gamma(Z^{0}\rightarrow\bar{f}f)$ decrease in the same manner by about $12$
to $15\%$. This decrease can be described by a universal function
$R(\varkappa)$ [see (56)–(60)] depending only on the energy-momentum of the
$Z^{0}$-boson and on the external-electromagnetic-field strength. This
circumstance becomes quite obvious if we consider that all known fermions
observed in $Z^{0}$-boson decays have masses at least an order of magnitude
smaller than the $Z^{0}$-boson mass. Therefore, it comes as no surprise that,
in the region of strong electromagnetic fields, the massless-fermion
approximation works well, also making it possible to calculate the dependence
of the total decay width $\Gamma_{Z}(\varkappa)$ of the $Z^{0}$-boson on the
external-field-strength parameter $\varkappa$ (9).
The results of the present calculations have revealed that, in just the same
way as in the case of the $W$-boson [23], the total decay width of the
$Z^{0}$-boson in a strong electromagnetic field depends nonmonotonically on
the strength parameter $\varkappa$, this dependence featuring the point of a
local minimum, $\Gamma_{Z}(\varkappa_{\rm min})=2,164$ GeV, at $\varkappa_{\rm
min}=0,445$. The graph in Fig. 11, which represents the behavior of
$\Gamma_{Z}(\varkappa)$ in the region $\varkappa\leq 1$, illustrates this
observation.
Figure 11: Minimal value of the total decay width of the $Z^{0}$-boson in an
external field.
For the sake of comparison, it can be indicated that the analogous $W$-boson-
decay process $W^{-}\rightarrow\ell\bar{\nu}_{\ell}$ in an external field is
also characterized by a nonmonotonic dependence of the decay width on the
parameter $\varkappa$ with the point of a local minimum at $\varkappa_{\rm
W}=0,6112$ [23]. Upon rescaling this value to the $Z^{0}$-boson energy scale,
it corresponds to $\varkappa=\varkappa_{\rm
W}\left(M_{W}/M_{Z}\right)^{3}=0.419$. Thus, the probabilities for the decays
of gauge bosons in an external field, $Z^{0}\rightarrow\bar{f}f$ and
$W^{-}\rightarrow\ell\bar{\nu}_{\ell}$, take their minimum values at nearly
the same field strength. This agreement is not accidental since the two
reactions in question proceed under similar kinematical conditions and since
the masses of the $W^{\pm}$ and $Z^{0}$ bosons are almost indistinguishable in
order of magnitude. It is also intriguing that, in the region around
$\varkappa\sim 1$, the presence or the absence of the initial-gauge-boson
electric charge is immaterial because the change in the decay width is due
primarily to the increase in the phase space of final particles.
In the region of superstrong fields ($\varkappa\geq 10$), a sizable
contribution to the total decay width of the $Z^{0}$-boson comes from the
$t$-quark-production process $Z^{0}\rightarrow\bar{t}t$, which is forbidden in
a vacuum by conservation laws. As the external-field strength increases, this
reaction becomes dominant in relation to other modes of $Z^{0}$-boson decay,
this being due to a very large mass of the $t$-quark. The external field
serves as some kind of a catalyst for super heavy-particle-production
processes, which cannot occur under ordinary conditions. In view of the
aforesaid, it would be of interest to discuss searches for new (as-yet-
undiscovered) hypothetical particles (supersymmetry and so on) among new
$Z^{0}$-boson decay modes that arise in an external electromagnetic field.
Although external-field strengths necessary for directly observing such
reactions are not yet available in experiments (see, for example, [21]), these
investigations are of course very interesting from the point of view of
fundamental principles of physical theory.
## APPENDIX
In the present study, we have employed the special mathematical functions
$\Phi(z),\Phi^{\prime}(z)$ and $\Phi_{1}(z)$ generically referred to as Airy
functions. The function $\Phi(z)$ is a particular solution to a second-order
linear differential equation for specific initial conditions:
$\Phi^{\prime\prime}(z)-z\Phi(z)=0,\quad\Phi(0)=\frac{\pi}{3^{2/3}\Gamma(2/3)},\quad\Phi^{\prime}(0)=-\frac{\pi}{3^{1/3}\Gamma(1/3)}.$
(A.1)
It has the well-known integral representation
$\Phi(z)=\int\limits_{0}^{\infty}\cos\left(zt+{t^{3}\over 3}\right)dt.$ (A.2)
Two other functions $\Phi^{\prime}(z)$ and $\Phi_{1}(z)$ can be expressed in
terms of this function as
$\Phi_{1}(z)=\int\limits_{z}^{\infty}\Phi(t)dt,\qquad\Phi^{\prime}(z)=\frac{d\Phi(z)}{dz}.$
(A.3)
At small values of the argument $z$ ($z\ll 1$), the Airy function $\Phi(z)$
can be calculated on the basis of the following expansion in a numerical
series:
$\Phi(z)={1\over 3^{2/3}}\sum_{n=0}^{\infty}\Gamma\biggl{(}{n+1\over
3}\biggr{)}\sin\biggl{(}{2\pi\over 3}+{2\pi n\over
3}\biggr{)}\frac{3^{n/3}}{n!}\cdot z^{n}.$ (A.4)
At large values of the argument ($z\gg 1$), the Airy function can be evaluated
by using the asymptotic expressions
$\Phi(z)={1\over 2z^{1/4}}\exp\left(-{2\over
3}z^{3/2}\right)\sum_{n=0}^{\infty}(-1)^{n}\frac{\Gamma(3n+1/2)}{(2n)!\
9^{n}}\cdot\frac{1}{z^{3n/2}},$ (A.5)
$\Phi(-z)={1\over z^{1/4}}\sum_{n=0}^{\infty}\sin\left({2\over
3}z^{3/2}+{\pi\over 4}-\frac{\pi n}{2}\right)\frac{\Gamma(3n+1/2)}{(2n)!\
9^{n}}\cdot\frac{1}{z^{3n/2}}.$ (A.6)
The properties of the Airy functions are well known and can be found in
mathematical handbooks (see, for example, [27]).
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|
arxiv-papers
| 2013-09-11T10:21:06 |
2024-09-04T02:49:50.703315
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Alexander V. Kurilin",
"submitter": "Alexander Kurilin",
"url": "https://arxiv.org/abs/1309.2780"
}
|
1309.2864
|
# Confined colloidal crystals in and out of equilibrium
A. Reinmüller Institut für Physik, Johannes-Gutenberg-Universität Mainz,
55128 Mainz, Germany E. C. Oğuz Institut für Theoretische Physik II: Weiche
Materie, Heinrich-Heine-Universität Düsseldorf, 40225 Düsseldorf, Germany R.
Messina Institut de Chimie, Physique et Matériaux (ICPM), Université de
Lorraine, 1 Bd. Arago, 57078 Metz - Cedex 3, France H. Löwen Institut für
Theoretische Physik II: Weiche Materie, Heinrich-Heine-Universität Düsseldorf,
40225 Düsseldorf, Germany H. J. Schöpe Institut für Physik, Johannes-
Gutenberg-Universität Mainz, 55128 Mainz, Germany T. Palberg Institut für
Physik, Johannes-Gutenberg-Universität Mainz, 55128 Mainz, Germany
###### Abstract
Recent studies on confined crystals of charged colloidal particles are
reviewed, both in equilibrium and out of equilibrium. We focus in particular
on direct comparisons of experiments (light scattering and microscopy) with
lattice sum calculations and computer simulations. In equilibrium we address
buckling and crystalline multilayering of charged systems in hard and soft
slit confinement. We discuss also recent crystalline structures obtained for
charged mixtures. Moreover we put forward possibilities to apply external
perturbations, in order to drive the system out of equilibrium. These include
electrolyte gradients as well as the application of shear and electric fields.
## I Introduction
Under strong confinement, as realized by system boundaries or laser-optical
external fields, the freezing transition of a colloidal suspension is not only
significantly shifted relative to its bulk freezing point Loewen_1994 but the
crystalline structure itself can differ much from its stable bulk phase
Lowen2012_Advances ; Schoepe2006_Langmuir ; Fontecha_2007 ; Fontecha2008_PRE ;
apolinario2006inhomogeneous ; brunner2002phase ; reichhardt2002novel ;
fortini2006 ; Qi2013_arXiv ; Rice2009_ChemPhysLett ; Erdal_JPCM . An example
of very strong confinement is given in a slit–geometry between two parallel
plates which results in a few crystalline layers if the plate distance is
small murray1990comparison ; Erdal_2009_EPL . The extreme limit is a quasi-
two-dimensional colloidal monolayer. Upon widening the confinement, the latter
can buckle into a bilayer pieranski1983thin ; Peeters_PRL_1999 ; Messina_2003
; Fontecha2005_JPCM ; Schoepe2006_Langmuir ; fortini2006 and subsequent
multilayered structures do emerge. Recent research has focused on the
structural details of the resulting multilayered crystal. Colloids are
excellent model systems to observe these structures on the particle scale
using optical microscopy as they can be confined between glass plates.
Moreover, various kinds of external fields can be used to bring the system out
of equilibrium Lowen2001 ; Villanova-Vidal2009_PhilosMag ; Schoepe2003_JCPM ;
Marques-Huesco2008_incoll resulting in novel phenomena not seen in
equilibrium or to study the dynamics of the crystallization process in detail.
The nature of the confinement can most easily be described by an external one-
body potential acting on the colloidal particles, for recent examples see
chakrabarti1995reentrant ; bechinger2001phase ; strepp2002phase ;
Resnick2003_JColloidInterfaceSci ; Brunner2004_PRL ;
apolinario2006inhomogeneous ; brunner2002phase ; reichhardt2002novel ;
mikhael2008archimedean ; schmiedeberg2008colloidal ; Schmiedeberg2007_EPJE .
One can therefore distinguish between hard and soft confinement: Hard
confinement implies a hard external potential, i.e. a potential which is
either zero or infinity describing just a steric volume exclusion of the
particles. Such an external interaction is athermal, i.e. if scaled with the
thermal energy, it is the same at different system temperatures which can
reduce the thermodynamic parameters of the total system. Soft confinement, on
the other hand, is described by a smooth external potential. An important
example are charged glass plates which act as a soft confinement to charged
colloids as described by an exponentially screened Coulomb interaction
hansen2000 ; Lowen2008_JCPM ; Oguz_2012 .
The goal of this minireview is twofold. First, we summarize some recent
progress in the experimental determination and theoretical prediction of
colloidal multilayers in various kinds of hard and soft confinement under
equilibrium conditions. As a result, the stable equilibrium structures reveal
a fascinating complexity being very sensitive to both particle interactions
and kind of confinement. Most of the structures can be understood by a simple
potential energy minimization, i.e. on purely energetic grounds. Second, we
focus on nonequilibrium situations which originate from the application of an
external field, like gravity, shear flow and an electrolyte gradient. We show
that, under appropriate solvent flow conditions, a seed can generate the
controlled growth of a colloidal monolayer. Moreover, we discuss how gravity
and electric fields can change crystallization. Last but not least, we
describe a complex self-organization process of colloidal clustering into a
swimming object in an electrolyte gradient. The physical reasons underlying
the details of the swimming process still need to be understood and clarified
in future studies.
The minireview is organized as follows: in chapter 2, we summarize all aspects
of confined colloidal crystal in equilibrium. Nonequilibrium effects as
induced by different external fields are subsequently discussed in chapter 3
and we conclude in chapter 4.
## II Phase behaviour of confined colloidal crystals in equilibrium
### II.1 Charged colloids in soft confinement
Regarding crystallization in slit geometry, hard confinement has been studied
for both hard particles and soft particles i.e. particles interacting with a
soft pair potential like a Yukawa interaction. More recently, the case of hard
particles in soft confinement has also been addressed Curk2012_PRE . Complex
cascades of multilayered crystals were predicted and confirmed in colloidal
experiments the details of which depend on the interparticle and particle–wall
interactions. Detailed comparisons between experiment and theory were
performed for hard and soft particles in hard confinement but for soft
particles in soft confinement such a direct comparison was done only recently
Oguz_2012 , which we describe in more detail in the following.
A flexible realization of soft sphere systems is provided by charged colloids
suspended in water. The screened Coulomb interaction can conveniently be
varied between the theoretical limits of hard spheres (approached at high salt
concentrations) and the One Component Plasma (OCP, approached at low salt
concentrations). To investigate charged colloidal spheres under deionized
conditions in confinement between parallel flat walls, Reinmüller et al.
realized a new experimental set–up to be used in combination with a commercial
scientific microscope Reinmueller2013_RevSciInst . An experimental set–up for
flexible and accurate investigations of such systems under confinement
requires a number of minimum properties: It must allow for fast and precise
adjustment of the confining geometry; it must facilitate efficient preparation
of colloidal suspensions under low salt concentrations at absence of unwanted
chemical gradients; it must provide an efficient homogenization mechanism as
well as long term stable experimental conditions. The measurement cell of the
setup consists of optically flat quartz substrates attached to piezo actuators
providing a fast and flexibly adjustable confining geometry. The local wall
separations and thus the confining geometry can be quantitatively controlled
in situ by interferometric measurements. Proper choice of materials guarantees
sufficient chemical inertia against contamination of any kind. Sample
preparation and homogenization is performed efficiently using a syringe pump
in combination with a conventional closed tubing system including a mixed bed
ion exchanger column. Combining optical microscopy, conductivity measurements
as well as interferometric measurements, the setup facilitates quantitative
control of particle interaction parameters and confining geometry as well as
the determination of the structural and dynamic properties of the confined
colloids.
Two mono-disperse aqueous suspensions of negatively charged polystyrene
spheres with diameters of 5.2$\mu$m and 2.6$\mu$m (batch nos PS/Q-F-B1036 and
PS-F-B233 by MicroParticles Berlin GmbH, Germany) were used to explore the
equilibrium phase diagram of charged colloidal spheres in aqueous suspensions
under spatial confinement in parallel plate geometry at low background salt
concentrations using the home-made setup described above. In order to suppress
gravity, the solvent mass density was matched to that of polystyrene by adding
20 vol% glycerol. During the experiments, the chemical properties of both
particle and substrate surfaces were kept constant. The equilibrium phase
diagram was determined in terms of the emerging crystal structures depending
on the dimensionless parameters the reduced area number density
$\eta=n_{A}d^{2}$ and the reduced inverse screening length $\lambda=\kappa d$,
where $n_{A}$ denotes the area particle number density, $\kappa$ is the
inverse screening length, and $d$ is the width of the confining slit. These
parameters could be determined by accurate measurements of the area number
density, the cell height, the salt concentration and the particle charge.
Varying the reduced area number density $\eta$ and the reduced inverse
screening length $\lambda$ in the experiment, the sequence of crystalline
structures shown in Figures 1 and 2 was observed. Basically these are
hexagonal closed packed crystals (fcc or hcp) oriented in different
crystallographic directions with respect to the cell walls. There are no prism
structures, transient structures (e.g. hcp like) or exotic structures, like
the Belgian waffle iron structure Erdal_2009_EPL , observable.
Figure 1: Examples of crystaline structures observed in experiments on charged
spheres in soft parallel plate confinement: $1\triangle$ (A), $2\square$ (B),
coexistence of $2\square$ and $2\triangle$ (C), $2\textrm{hcp}\perp$ (central
large domain) in coexistence with $2\triangle$ (grain at the right margin)(D),
$3\triangle$ with different appearances of fcc(111) and hcp(001) faces (E) and
Moiré rosettes in coexistence with $2\textrm{hcp}\perp$ and $2\triangle$ (F).
(Field of view: $280\times 210\mu\textrm{m}^{2}$.) From Ref. Oguz_2012 .
The experimentally observed equilibrium phase behavior of highly charged
colloidal crystals in parallel soft confinement was found to be in very good
agreement with the theoretical ground state phase diagrams of expponentially
screened point-like Coulomb (Yukawa) particles calculated by lattice sum
minimizations. The particles were exposed to a soft confining hyperbolic
cosine potential as described by the linear screening theory Andelman . The
phase diagrams display the same simple multilayer phase sequence with
increasing system density. The agreement is documented in Figure 2 which
contains a comparison between experiment and theory in the two–dimensional
plane spanned by $\eta$ and $\lambda$ for the $5.2\mu\textrm{m}$ size
colloids.
Figure 2: Comparison of the theoretical prediction (dashed and full lines) for
the multilayer stability phase diagram at zero temperature to the experimental
data (symbols) taken from the phase diagram for PS $5.2\mu\textrm{m}$
particles in Oguz_2012 . We use $\kappa a=6$ and a particle/wall surface
charge ratio of $\nu=0.4$, where $a$ is the particle radius. From Ref.
Oguz_2012 .
### II.2 Hard slit confinement
One of the simplest situation of confined colloids are sterically-stabilized
particles (”hard spheres”) confined between parallel hard walls. In this slit
geometry the packing fraction and the reduced plate distance $L/\sigma$ (with
$\sigma$ denoting the hard sphere diameter) are the only thermodynamic
parameters Schmidt1 ; Schmidt2 . The two-dimensional limit of extreme
confinement, $L\to\sigma$, gives a system of hard disks. In this case, the
close-packed situation is a hexagonal (or triangular) two-dimensional crystal.
For finite packing fractions, the nature of the freezing transition has
attracted large attention in the last years Bernard2011 ; Peng10 ;
Engel2013_PRE . For slightly larger plate distances, the spheres deviate a bit
from their ideal mid-plane positions. In this case, the perturbative
correction to the exact partition function can be calculated analytically
Franosch2012 . The cascade of close-packing sequences as a function of
$L/\sigma$ is nontrivial once the distance exceeds that of two intersecting
hexagonal layers, see ivlev2012complex for a review. The phase diagram as
determined from free-energy calculations was first investigated in fortini2006
, which reveals the thermodynamic stability of close-packed structures such as
prisms. More recently the best packed configuration was found by systematic
numerical minimization routines Oguz_2012_PRL . As a result, periodic adaptive
prismatic structures which are composed of alternating prisms of spheres are
best packed. The internal structure of these prisms adapts to the slit height
which results in close packings for a range of plate separations, just above
the distance where three intersecting square layers fit exactly between the
plates. The adaptive prism phases were also observed in real-space experiments
on confined sterically stabilized colloids and in Monte Carlo simulations at
finite pressure Oguz_2012_PRL .
### II.3 Cylindrical confinement
If particles are confined in cylindrical tubes, there is no sharp freezing
transition since phase transitions in general are rounded in one dimension by
fluctuations Wilms . However, the ground state at zero temperature can be a
periodic crystal. It has been shown, that hard spheres confined to a hollow
hard cylinder show complex ordering structures with helical order Pickett2000
. The ground state structures were generalized recently to Yukawa systems in a
hard cylindrical pore Oguz_2011 . As a function of screening strength and
particle density, the equilibrium phase diagram was found to exhibit a cascade
of stable crystals with both helical and non-helical structures.
### II.4 Monolayers of oppositely charged colloids
Oppositely charged colloids form stable crystals even at zero pressure due to
the mutual attractive Coulomb forces, for a recent review see ivlev2012complex
. A model of oppositely charged hard spheres with different diameters is a
standard description for the basic interactions. In three dimensions the
ground state at zero pressure is textbook knowledge evans1966crystal41to43 .
There are three basic equimolar structures whose stability depends on the
ratio of the ion radii involving the cesium–chloride, sodium–chloride, and
zincblende structures. These (and some more) structures are not only realized
for molecular salts but also for mesoscopic colloids leunissen2005 ;
hynninen2006 ; hynninen2006cuau . In Refs. Assoud2010_EPL , the same problem
was solved for two-dimensional monolayers of oppositely charged hard spheres.
Among the stable structures are square, triangular and rhombic crystals as
well as “empty” crystals made up of dipoles and chains, which have a vanishing
number density. The square structure was already observed for charged
granulates Kaufman2009 . As an example, Figure 3 shows the cascade of zero-
pressure ground states as a function of the diameter ratio for an interface
set–up (interfacial model) where the centers of mass of the two species are
all in a common plane. The associated ground state energy per particle is
plotted versus the diameter ratio.
Figure 3: Stable structures of oppositely charged spheres versus their size
asymmetry $R_{B}/R_{A}$ in the interface model, where all sphere centers fall
on the same plane: a) side view, b) (scaled) energy per ion ($q$ is the
particle charge). The discontinuous transition is indicated by a solid bar.
Continuous transitions are denoted by a broken bar. Unit cells of the
corresponding stable phases are shown, where the big (small) spheres have a
radius $R_{A}$ ($R_{B}$). From Ref. Assoud2010_EPL .
## III Phase behaviour of confined colloidal crystals out of equilibrium
### III.1 Influence of gravity
The influence of gravity on the crystal formation in a two layer system of
charged colloidal particles in confinement was explored experimentally using
commercial polystyrene spheres of 5.2 $\mu$m in diameter in fully deionized
water Moire . Here, the colloidal bilayer crystal assembly occurs under
influence of both gravity and electrostatics. During a first stage of
settling, bilayers form via a line pattern, which quickly gives way to a
polycrystalline bilayer structure of AB–stacking of two triangular monolayers.
Due to gravity, this bilayer structure is a unidirectional distorted crystal:
The lattice constant in vertical direction is noticeable smaller than in
lateral direction. In addition to the triangular lattice, two extraordinary
transient structures were identified. Each is made from two perfect triangular
layers which are rotated by certain angles with respect to each other. The
intergrowth of the conventionally stacked two layer crystals results in
$1\times 1R\alpha$ super structures. Two specific rotation angles
$\alpha=27.8^{\circ}$ and $\alpha=38.2^{\circ}$ dominate in the observed two
layer system. These superstructures show fascinating triangular Moiré patterns
of rosettes. The morphological function of these structures can be
characterized as an “extended grain boundary” known from atomic systems. A
similar coherent intergrowth of crystal domains in the bulk are observed e.g.
in rare earth boride carbides. Due to their complex structure, the observed
crystalline arrangements display extraordinary scattering patterns showing a
twelve-fold symmetry. Measurements and model calculations demonstrate the high
potential of these extraordinary structures for photonic applications.
### III.2 Crystal formation in parallel confinement under the influence of
shearing fields
Regarding crystallization in confined geometry, heterogeneous nucleation at
the cell wall and the influence of shearing fields are of substantial
importance. Colloidal crystals are easily shear molten by gentle mechanical
agitation. After cessation of shear, they readily crystallize via nucleation
and growth. Applying randomly directed shear before crystallization bcc
crystals nucleating at the wall are oriented with their (110)-bcc plane
parallel to the wall, but do not show any preferred orientation of the
(111)–direction Engelbrecht2011_SoftMatter . The crystal growth in vertical
direction was observed to be reaction controlled with a constant growth
velocity and as function of meta-stability the growth velocities follow a
Wilson-Frenkel-law. The resulting polycrystalline material consist of columnar
crystals randomly orientated in the crystallographic direction parallel to the
cell wall.
Analyzing the heterogeneous nucleation rate densities on the wall as function
of metastability, a transition from a nearly time independent nucleation rate
density at the phase boundary to a peaked transient nucleation at higher
undercooling $\Delta\mu$ was identified. The heterogeneous nucleation barrier
height stays first constant and drops down to zero as the metastability
increases: In charged sphere model systems, a sharp wetting transition
separates a regime at moderate metastability where hemispherical cap nuclei
are formed from a strong meta-stability regime, where the crystal phase
completely wets the wall.
In suspensions molten by unidirectional shear flow this, transition could not
be observed. Unidirectional shear leads to shear induced oriented precursor
crystals in form of hexagonal layers at the cell wall ((111)–fcc parallel to
the wall) and further quasi-epitaxial growth of bcc crystals on this precursor
crystal Wette2009_JCPM ; Palberg2012_JCP . The precursor acts as structured
substrate reducing the nucleation barrier height to a minimum. The details of
this first step are still not well understood. Due to the unidirectional
shear, the bcc crystals are uniformly oriented ((110) plane parallel to the
wall and (111)–direction in direction of shear flow) leading to a twinned bcc
morphology. Like in the first case, vertical growth is reaction controlled
with a constant growth velocity. In lateral dimension, the average domain size
was observed to increase significantly during vertical growth following the
typical power law for coarsening with an exponent of 0.5. Lateral coarsening
ceases when the wall crystals from the opposing walls intersect.
These investigations and others Messina_2006 show impressively that the
heterogeneous nucleation process on the wall, the crystallization process
afterwards and the crystalline microstructure is highly modified by shear in
confined geometry. Finally, rigid colloidal spheres rotate in shear flow.
Flexible neutral polymer chains, on the other hand, tumble He_2010 . If
combined with other periodic field there are interesting dynamical modes of
polymer propagation He_2011 . The behaviour is expected to be even more
complex for charged polymers Allahyarov_2011 .
### III.3 Crystal and microswimmer formation in parallel confinement under
the influence of electrostatic and flow fields
#### III.3.1 Crystal nucleation
The crystallization process in equilibrium two–dimensional colloidal fluids in
confinement induced by local electrolyte gradients was studied experimentally
using negatively charged polystyrene spheres (diameter $2a=5.2\mu\textrm{m}$)
dispersed in water Reinmueller2012 . The suspension was prepared under
deionized conditions in confined parallel plate geometry obtaining a fluid
monolayer system in thermodynamic equilibrium. After introducing small
fragments of ion exchange resin into the confinement cell, a strong
interaction between the colloidal spheres and the fragments could be observed.
In particular for the case of cation exchange resins (CIEX) a fluid flow leads
to particle accumulation at the CIEX-fragment followed by colloidal
crystallization. The electrolyte gradient originated by the CIEX-fragment
caused an electro-osmotic solvent flow along the anionically charged substrate
surface towards the fragment leading to a colloidal particle transport. Once
the local particle number density exceeded the salt concentration dependent
freezing density, heterogeneous crystallization on the surface of the CIEX-
fragment occurred. Although the underlying mechanism generating the flow field
is not fully understood, the phenomenology of the crystal formation was
observed to be well reproducible. Starting from the thermodynamically stable
colloidal fluid multi-domain triangular monolayer crystals were formed with
the fragment in the centre, as for experimental snapshots see Figure 4. The
domain number is highly correlated with the shape of the fragment acting as
seed particle. Analyzing the crystal growth, an exponentially decreasing
crystal growth velocity was observed and the final crystallite size is
determined by the strength and range of the electrolyte gradients. The same
situation was simulated using Brownian dynamics computer simulations in two
dimensions and good agreement with the experiment was obtained.
Figure 4: Growth of a three–domain monolayer crystal on a seed at
$t=0\textrm{s}$ (a), $t=100\textrm{s}$ (b), $t=200\textrm{s}$ (c),
$t=300\textrm{s}$ (d), $t=400\textrm{s}$ (e), and $t=500\textrm{s}$ (f); field
of view: $360\times 290\mu\textrm{m}^{2}$. From Ref. Reinmueller2012 .
Simulation and experiment indicate that there is no nucleation barrier in this
special crystallization mechanism and that the number of domains depends
neither on thermodynamics nor on nucleation kinetics. The seed geometry
determines the microstructure of the formed crystal, see simulation snapshots
in Figure 5 for two different seed geometries. In principle these
investigations offer the possibility of designing and manufacturing extended
thin multidomain colloidal crystal microstructures controlling the positions,
shape and size of individual crystals.
Figure 5: Simulation snapshots of a three–domain (a) and a mono–domain crystal
(b) each at $t=500\tau$ with $\tau$ being a typical Brownian time scale. The
green dots in both illustrations indicate the origin of the system. The
particle flow into the polygonal (a) and spherical (b) area is prohibited. The
red particles are crystalline, while the black ones are not. Field of view:
$250\times 250/\kappa^{2}$ where $\kappa$ is the inverse screening length.
From Ref. Reinmueller2012 .
#### III.3.2 Complex microswimmer formation
Repeating the experiment with stronger electro–osmotic solvent flow and larger
plate to plate distance ( $\sim 50\mu\textrm{m}$) the formation of self-
organized complex microswimmers could be observed Reinmueller2013_Langmuir .
The electrolyte gradient inducing a diffusio-osmotic solvent flow towards the
CIEX fragment now originates a convection roll at the CIEX fragment. Particles
transported to the fragment were lifted up at the CIEX and fall down afterward
very close to the fragment. Due to the convection the particles did not
crystallize, but still a strong accumulation of particles at the fragment was
observed. In some cases the particle concentration within the convection cell
was found to be highly asymmetric (rotational symmetry, symmetry axis in the
centre of the CIEX) which might be generated by either a shape anisotropy of
the CIEX fragment or by a highly asymmetric particle concentration of the
colloidal fluid surrounding the CIEX fragment.
If the particle concentration within convection cell was highly asymmetric the
CIEX fragment and the accumulated particles started to move in the direction
of smallest colloid concentration. A possible propulsion mechanism related to
the nonuniform colloid concentration (including colloid counterions) might be
direct osmotic propulsion. The details of the driving mechanism of this
cooperative motion are not yet fully understood. Nevertheless the propulsion
was found to be very robust and occurred with varying effectiveness over a
large range of experimental parameters. The self propulsion mechanism was
self-stabilizing over many minutes and the complex micro swimmers moved with
typical velocities of $1\mu\textrm{m}/\textrm{s}-3\mu\textrm{m}/\textrm{s}$.
At sufficient speed, further loading from the side became negligible and
oncoming particles fall down behind the CIEX fragment after they were lifted
up by convective flow. During this stage of constant trail thickness, the
complexes cover path lengths of many hundreds of micrometers.
### III.4 Oppositely charged colloidal suspensions in electric fields: laning
When oppositely charged colloids are exposed to an external electric field,
laning of like-charge particles occurs, for a review see ivlev2012complex ;
Sutterlin09 ; Dzubiella02 ; lowen2010instabilities . In three dimensions, by
experiment and computer simulation, laning was shown to be continuous for
small electric fields and a quantitative comparison between experiment and
simulation could be achieved vissers2011 . It is important to note that
hydrodynamic interactions between driven particle are of minor importance when
charged particle are moving in an electric field Ajdari_Long_EPJE_2001 ;
Rex_EPJE_2008 . Therefore they were not found to be crucial for the actual
modelling.
## IV Conclusions
In conclusion, we have reviewed briefly recent developments in preparing,
characterizing and predicting confined colloidal crystals. In equilibrium, a
huge variety of equilibrium structures becomes stable depending sensitively on
the particle interaction and on the nature of the confinement. These
structures may serve as building blocks for filters and sieves as well as for
photonic crystals. A coagulated two-dimensional crystal has controlled void
space between the particles, which leads to tunable porosity. We have also
shown that nonequilibrium effects generated by shear flow, electrolyte
gradients and electric fields may generate further structures which are widely
different from those stable in equilibrium. These external fields can also be
used to steer the crystallization process itself and control the defect
structure and distribution of the emerging crystal. One of the most striking
consequence was the self-organized formation of a microswimming system which
opens the way to possible applications as controlled cargo transport via
active colloidal suspensions.
Future research activities should focus on microscopic theories for
nonequilibrium processes in order to predict nonequilibrium effects of
crystallization. Dynamical density functional theory is a promising approach
in this respect seed ; Espanol2009 ; Advances . On the experimental side,
systems with depletion attractions and binary and ternary mixtures should be
explored more. In these cases, new segregation and aggregation phenomena are
expected which may lead to further applications. Therefore, though the basic
mechanisms are understood by now, a flourishing future of crystallization in
confined colloidal suspensions is still lying ahead.
## Acknowledgements
This work was supported by the DFG within SFB TR6 (project D1).
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|
arxiv-papers
| 2013-09-11T15:40:55 |
2024-09-04T02:49:50.718922
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "A. Reinm\\\"uller, E. C. O\\u{g}uz, R. Messina, H. L\\\"owen, H. J.\n Sch\\\"ope, and T. Palberg",
"submitter": "Erdal Celal O\\u{g}uz",
"url": "https://arxiv.org/abs/1309.2864"
}
|
1309.2881
|
# Hidden $T$-Linear Scattering Rate in Ba0.6K0.4Fe2As2 Revealed by Optical
Spectroscopy
Y. M. Dai Beijing National Laboratory for Condensed Matter Physics, Institute
of Physics, Chinese Academy of Sciences, P.O. Box 603, Beijing 100190, China
LPEM, ESPCI-ParisTech, CNRS, UPMC, 10 rue Vauquelin, F-75231 Paris Cedex 5,
France Condensed Matter Physics and Materials Science Department, Brookhaven
National Laboratory, Upton, New York 11973, USA B. Xu Beijing National
Laboratory for Condensed Matter Physics, Institute of Physics, Chinese Academy
of Sciences, P.O. Box 603, Beijing 100190, China LPEM, ESPCI-ParisTech, CNRS,
UPMC, 10 rue Vauquelin, F-75231 Paris Cedex 5, France B. Shen H. Xiao
Beijing National Laboratory for Condensed Matter Physics, Institute of
Physics, Chinese Academy of Sciences, P.O. Box 603, Beijing 100190, China H.
H. Wen Beijing National Laboratory for Condensed Matter Physics, Institute of
Physics, Chinese Academy of Sciences, P.O. Box 603, Beijing 100190, China
National Laboratory of Solid State Microstructures and Department of Physics,
Nanjing University, Nanjing 210093, China X. G. Qiu Beijing National
Laboratory for Condensed Matter Physics, Institute of Physics, Chinese Academy
of Sciences, P.O. Box 603, Beijing 100190, China C. C. Homes [email protected]
Condensed Matter Physics and Materials Science Department, Brookhaven National
Laboratory, Upton, New York 11973, USA R. P. S. M. Lobo [email protected] LPEM,
ESPCI-ParisTech, CNRS, UPMC, 10 rue Vauquelin, F-75231 Paris Cedex 5, France
(August 27, 2024)
###### Abstract
The optical properties of Ba0.6K0.4Fe2As2 have been determined in the normal
state for a number of temperatures over a wide frequency range. Two Drude
terms, representing two groups of carriers with different scattering rates
($1/\tau$), well describe the real part of the optical conductivity,
$\sigma_{1}(\omega)$. A “broad” Drude component results in an incoherent
background with a $T$-independent $1/\tau_{b}$, while a “narrow” Drude
component reveals a $T$-linear $1/\tau_{n}$ resulting in a resistivity
$\rho_{n}\equiv 1/\sigma_{1n}(\omega\rightarrow 0)$ also linear in
temperature. An arctan($T$) low-frequency spectral weight is also a strong
evidence for a $T$-linear 1/$\tau$. Comparison to other materials with similar
behavior suggests that the $T$-linear $1/\tau_{n}$ and $\rho_{n}$ in
Ba0.6K0.4Fe2As2 originate from scattering from spin fluctuations and hence
that an antiferromagnetic quantum critical point is likely to exist in the
superconducting dome.
###### pacs:
72.15.-v, 74.70.-b, 78.30.-j
Over the last several decades, it has been observed that the electrical
resistivity $\rho$ of some strongly-correlated materials increases linearly
with temperature ($T$-linear $\rho$), particularly in the vicinity of an
antiferromagnetic quantum critical point (QCP), a striking deviation from
Landau’s Fermi-liquid description of metals. This anomalous $T$-linear $\rho$,
extensively studied in the high-$T_{c}$ cuprate superconductors Boebinger
(2009); Cooper _et al._ (2009); Daou _et al._ (2009); Jin _et al._ (2011),
organic Bechgaard salts Doiron-Leyraud _et al._ (2009a, b), as well as heavy-
fermion metals Malinowski _et al._ (2005); Tanatar _et al._ (2007);
Trovarelli _et al._ (2000); Gegenwart _et al._ (2002), may be intimately
related to the emergence of superconductivity Jin _et al._ (2011); Taillefer
(2010); Doiron-Leyraud _et al._ (2009a). It is generally believed that in
proximity to the antiferromagnetic QCP, spin fluctuations are so strong that
the scattering process of quasiparticles is severely modified, inducing non-
Fermi-liquid behavior. Such spin fluctuations may be responsible for the
pairing of electrons in high-$T_{c}$ superconductors Monthoux _et al._
(2007); Moriya and Ueda (2000); Mazin _et al._ (2008).
The presence of the $T$-linear $\rho$ and a QCP in the newly discovered iron-
based superconductors is highly desired since superconductivity arises in the
vicinity of the spin-density-wave (SDW) instability Kamihara _et al._ (2008);
Rotter _et al._ (2008). Up to this point, $T$-linear $\rho$ has been observed
by transport, arousing considerable effort to describe it as evidence of
possible quantum criticality in iron-pnictides, especially in the “122” family
Xu _et al._ (2008); Gooch _et al._ (2009); Shen _et al._ (2011); Maiwald
_et al._ (2012); Hashimoto _et al._ (2012); Kasahara _et al._ (2010);
Doiron-Leyraud _et al._ (2009a); Shibauchi _et al._ (2013). However, unlike
the high-$T_{c}$ cuprates, iron-pnictides fall into the category of multi-band
materials Singh and Du (2008); Ding _et al._ (2008). Up to five Fe-3$d$ bands
crossing the Fermi level contribute to the Fermi surface, leading to the
presence of three hole-like Fermi pockets at $\Gamma$-point and two electron-
like pockets at the corners of Brillouin zone. The scattering rate $1/\tau$
and the response of quasiparticles to external electrical field may vary
considerably in different Fermi pockets Fang _et al._ (2009); Shen _et al._
(2011); Muschler _et al._ (2009); Rullier-Albenque _et al._ (2010). As a
result, the transport properties become extremely complicated in such a system
and the question of whether a $T$-linear $\rho$ in iron-pnictides originates
from multi-band effects or the presence of a QCP makes it inadequate to
investigate the transport properties alone in search of a non-Fermi-liquid
behavior and evidence for possible QCP Rullier-Albenque _et al._ (2010,
2009).
Although many optical studies on iron-pnictides have been reported Li _et
al._ (2008); Charnukha _et al._ (2011a, b), the above issue has never been
touched due to the absence of detailed $T$ dependent optical data. In this
Letter we address this issue by studying the detailed $T$ dependence of the
optical conductivity and low-frequency spectral weight in Ba0.6K0.4Fe2As2. The
low-frequency optical conductivity is described by two Drude terms: a broad
Drude with a large $1/\tau_{b}$ that is basically $T$-independent alongside a
narrow one with a small $T$-linear $1/\tau_{n}$ that reveals a $\rho_{n}\equiv
1/\sigma_{1n}(\omega\rightarrow 0)\propto T$. The low-frequency spectral
weight increases with cooling, following an $\arctan(T)$ dependence, which is
demonstrated to be a clear signature of $T$-linear $1/\tau$. Comparison with
similar behavior found in other materials attributes the $T$-linear
$1/\tau_{n}$ and $\rho_{n}$ in Ba0.6K0.4Fe2As2 to spin fluctuation scattering
and the presence of a QCP in the superconducting dome.
High quality Ba0.6K0.4Fe2As2 single crystals were grown using a self-flux
method Shen _et al._ (2011). The inset of Fig. 1 shows the DC resistivity of
Ba0.6K0.4Fe2As2 as a function of temperature.
Figure 1: (color online) The thick solid lines in the main panel are ab-plane
$\sigma_{1}(\omega)$ of Ba0.6K0.4Fe2As2 at different temperatures in the
normal state. The smooth lines through the data are fits with Drude-Lorentz
model. The inset shows the DC resistivity as a function of temperature from
transport measurement (solid curve) and the values derived from the zero
frequency extrapolation of $\sigma_{1}(\omega)$ (solid circles).
The $\rho(T)$ curve is characterized by a sharp superconducting transition at
$T_{c}\simeq 39$ K and a tendency to saturation at room temperature. A visible
change in slope occurs at about 175 K.
The near-normal incident reflectance, $R(\omega)$, has been measured for light
polarized in the a-b plane using FTIR spectrometers and an in situ evaporation
technique Homes _et al._ (1993). Data from $\simeq 20-12\,000$ cm-1 were
collected at 18 different temperatures from 5 to 300 K on a freshly cleaved
surface. The visible and UV range ($10\,000-55\,000$ cm-1) $R(\omega)$ was
measured at room temperature with an AvaSpec-2048$\times 14$ fiber optic
spectrometer. The real part of the complex optical conductivity,
$\sigma_{1}(\omega)$, is determined from $R(\omega)$ via Kramers-Kronig
analysis. A Hagen-Rubens form ($R=1-A\sqrt{\omega}$) is used for the low-
frequency extrapolation. At high frequencies, $R(\omega)$ is assumed to be
constant to 40 eV, above which a free-electron response ($\omega^{-4}$) is
used.
The main panel of Fig. 1 shows $\sigma_{1}(\omega)$ at 4 selected temperatures
in the normal state (thick solid lines); all the spectra exhibit the well-
known Drude-like metallic response. In order to quantitatively analyze the
optical data, we fit $\sigma_{1}(\omega)$ to the Drude-Lorentz model,
$\sigma_{1}(\omega)=\frac{2\pi}{Z_{0}}\\!\left[\sum_{k}\frac{\Omega^{2}_{p,k}}{\tau_{k}(\omega^{2}+\tau_{k}^{-2})}\\!+\\!\sum_{j}\frac{\gamma_{j}\omega^{2}\Omega_{j}^{2}}{(\omega_{j}^{2}-\omega^{2})^{2}+\gamma_{j}^{2}\omega^{2}}\right]$
(1)
where $Z_{0}$ is the vacuum impedance. The first term describes a sum of free-
carrier Drude responses, each characterized by a plasma frequency
$\Omega_{p}^{2}=4\pi ne^{2}/m^{\ast}$, where $n$ is a carrier concentration
and $m^{\ast}$ is an effective mass, and a scattering rate $1/\tau$. The
second term corresponds to a sum of Lorentz oscillators characterized by a
resonance frequency ($\omega_{j}$), a line width ($\gamma_{j}$) and an
oscillator strength ($\Omega_{j}$). This Drude-Lorentz model is also used to
determine the dc properties of the system Romero _et al._ (1992).
As shown in Fig. 2, $\sigma_{1}(\omega)$ at 150 K is described by a broad
Drude with a large scattering rate $1/\tau_{b}\approx 936$ cm-1, and a narrow
Drude with a small scattering rate $1/\tau_{n}\approx 158$ cm-1 and an
overdamped Lorentzian term.
Figure 2: (color online) The black solid curve is the measured
$\sigma_{1}(\omega)$ of Ba0.6K0.4Fe2As2 at 150 K. The red dashed line through
the data is the fit which is decomposed into a broad Drude (blue dotted line),
a narrow Drude (pink dotted-dashed line) and a Lorentz (orange short-dashed
line) term.
The linear superposition of these three components gives very good description
to $\sigma_{1}(\omega)$ up to $10\,000$ cm-1 at all the measured temperatures
in the normal state. Fits for other temperatures are selectively shown in the
main panel of Fig. 1 as smooth thin lines through the corresponding data. The
inset of Fig. 1 compares the optical estimate for the DC resistivity
$\rho\equiv 1/\sigma_{1}(\omega\rightarrow 0)$ (solid circles) to the
transport measurements (solid line).
The two-Drude fit indicates the existence of two groups of carriers with
different $1/\tau$’s in Ba0.6K0.4Fe2As2, which was first pointed out by Wu _et
al._ in various iron-pnictides Wu _et al._ (2010). The disparity of the
$1/\tau$’s in different bands is also supported by both theoretical
calculations Kemper _et al._ (2011) and measurements using other techniques
on similar materials Fang _et al._ (2009); Shen _et al._ (2011); Rullier-
Albenque _et al._ (2010); Muschler _et al._ (2009). Tu et al. suggest that
it is appropriate to describe the broad Drude term as bound excitations Tu
_et al._ (2010), because the mean free path $l=v_{F}\tau$ ($v_{F}$ is the
Fermi velocity) associated with the broad Drude is less than the shortest
interatomic spacing, violating the Mott-Ioffe-Regel condition Gurvitch (1981).
In Ba0.6K0.4Fe2As2, the average Fermi velocities of the electron and hole
pockets are estimated to be $v^{e}_{F}\simeq 0.40$ eVÅ and $v^{h}_{F}\simeq
0.36$ eVÅ Ding _et al._ (2011), respectively. Furthermore, it is reported
that in iron-pnictides holes have a larger $1/\tau$ than electrons Fang _et
al._ (2009); Shen _et al._ (2011); Rullier-Albenque _et al._ (2010);
Muschler _et al._ (2009). If we associate the broad Drude component
($1/\tau_{b}\simeq 936$ cm-1) with the hole pockets, a mean free path of
$l_{h}\simeq 3$ Å is obtained. This value is close to the lattice parameter
$a\simeq 4$ Å of the 122 family compounds and probably too small for coherent
transport. Since the broad Drude only produces an incoherent, $T$-independent,
background contribution to the total $\sigma_{1}(\omega)$, the nature of the
broad Drude component (whether or not bound excitations) does not affect our
analysis on the temperature dependence of $\sigma_{1}(\omega)$ and low-
frequency spectral weight.
Figure 3: (color online) The $T$ dependence of (a) the plasma frequency
$\Omega_{p}$, (b) the scattering rate $1/\tau$, (c) the contribution to DC
conductivity $\sigma_{1}(\omega\rightarrow 0)$, and (d) the equivalent
resistivity $\rho$ for the narrow and broad Drude components, respectively.
The dashed lines in panels (b) and (d) are linear fits to $1/\tau_{n}$ and
$\rho_{n}$, respectively.
The temperature dependence of the two Drude components provides information
about the nature of the two different types of carriers in this material.
Figure 3 shows the $T$ dependence of the Drude parameters from our fits. The
subscripts n and b stand for the narrow and broad Drude terms, respectively.
Fig. 3(a) shows the $T$ dependence of the plasma frequencies of the two Drude
terms. Upon cooling, $\Omega_{p}$ is roughly a constant for each of the two
components, indicating that the band structure and $n/m^{\ast}$ does not
change with temperature, in agreement with a previous work Wu _et al._
(2010). Fig. 3(b) portrays the $T$ dependence of the scattering rate of the
two Drude components, where $1/\tau_{b}$ is basically $T$-independent while
$1/\tau_{n}\propto T$; the black dashed line denotes a linear fit. Fig. 3(c)
displays the contribution of the two groups of carries to the DC conductivity.
As the broad Drude ($\sigma_{b}$) is $T$-independent, the temperature
dependence of the total DC conductivity arises out of the narrow Drude band
($\sigma_{n}$). The $\rho(T)$ curve of Ba0.6K0.4Fe2As2, shown in the inset of
Fig. 1, exhibits a tendency to saturation at room temperature, and a change of
slope can be seen at about 175 K. This behavior can be explained by the
different $T$ dependence of the two Drude bands, which can be considered as a
parallel-circuit Wiesmann _et al._ (1977): $\sigma=\sigma_{n}+\sigma_{b}$. A
crossover-region, where $\sigma_{n}\simeq\sigma_{b}$, is found in Fig. 3(c) at
$\simeq 175$ K. Below this temperature, $\sigma_{n}>\sigma_{b}$, so the total
DC conductivity is dominated by $\sigma_{n}$ which exhibits strong temperature
dependence. As a result, below 175 K, $\rho(T)$ decreases quickly with
decreasing temperature. Above 175 K, $\sigma_{n}<\sigma_{b}$, and the total DC
conductivity is dominated by $\sigma_{b}$ which shows no temperature
dependence. Hence, above 175 K, the growth of the DC resistivity slows with
heating, resulting in the change of slope and the tendency to saturation in
$\rho(T)$. Similar conclusions were obtained from investigations of the Hall
effect Rullier-Albenque _et al._ (2010, 2009) and theoretical calculations
Golubov _et al._ (2011). In the $\omega\rightarrow 0$ limit, the inverse of
$\sigma_{1}(0)$ yields the resistivity from the two Drude components, as shown
in Fig. 3(d). The resistivity of the broad Drude remains a constant at all
measured temperatures, while a $T$-linear $\rho_{n}$ is revealed for the
narrow Drude component. This is in accord with transport measurements on hole-
doped 122 compounds Shen _et al._ (2011); Maiwald _et al._ (2012); Gooch
_et al._ (2009), where $T$-linear $\rho$ was observed at low temperatures in
optimally-doped materials. The $T$-linear $\rho$ is only found at low
temperatures as this is the region dominated by the narrow Drude component.
Further evidence for $T$-linear $1/\tau$ lies in the temperature dependence of
the low-frequency spectral weight. The spectral weight is defined as
$W_{0}^{\omega_{c}}=\int_{0}^{\omega_{c}}\sigma_{1}(\omega)d\omega,$ (2)
where $\omega_{c}$ is a cutoff frequency. In a Drude metal the scattering rate
decreases upon cooling, producing a narrowing of the Drude response and an
increase of the DC conductivity, resulting in a transfer of spectral weight
from high to low frequencies and an increase in $W_{0}^{\omega_{c}}$. To
quantitatively analyze the $T$ dependence of the low-frequency spectral
weight, we adopt one Drude optical conductivity [see Eq. (1)] into Eq. (2), to
obtain the spectral weight as a function of $1/\tau$
$W_{0}^{\omega_{c}}(1/\tau)=\frac{2\pi\Omega_{p}^{2}}{Z_{0}}\arctan\left(\omega_{c}\tau\right).$
(3)
In the case of $1/\tau\propto T$, Eq. (3) can be simplified as
$W_{0}^{\omega_{c}}(T)=a_{1}\arctan\left(\frac{a_{2}}{T}\right),$ (4)
where $a_{1}={2\pi\Omega_{p}^{2}}/Z_{0}$, and $a_{2}\propto\omega_{c}$; both
are $T$-independent parameters. Considering the spectral weight arising from
the incoherent part and inter-band transition (Lorentz), which are both
$T$-independent, we introduce the third $T$-independent parameter $a_{0}$ into
Eq. (4). Finally, the low-frequency spectral weight as a function of
temperature for $1/\tau\propto T$ is written as
$W_{0}^{\omega_{c}}(T)=a_{0}+a_{1}\arctan(\frac{a_{2}}{T}).$ (5)
$W_{0}^{\omega_{c}}$ can be easily determined by integrating the measured
$\sigma_{1}(\omega)$. The open symbols in Fig. 4 denote $W_{0}^{\omega_{c}}$
with different $\omega_{c}$’s: 150 (squares), 200 (triangles) and 250 cm-1
(diamonds).
Figure 4: (color online) Temperature dependence of the low-frequency spectral
weight $W_{0}^{\omega_{c}}$ with different $\omega_{c}$’s: 150 (squares), 200
(triangles) and 250 cm-1 (diamonds). The solid lines through the data
represent the expected temperature dependence for $1/\tau\propto T$ at each
$\omega_{c}$.
As expected in a metallic system, for all the $\omega_{c}$’s,
$W_{0}^{\omega_{c}}$ increases with decreasing $T$. The solid curves through
the data are the least square fits using Eq. (5). The excellent agreement
indicates that a $T$-linear $1/\tau$ dominates the low frequency
$\sigma_{1}(\omega)$ over a very broad temperature range.
We now discuss the possible origin of the $T$-linear $1/\tau_{n}$ and
$\rho_{n}$. A $T$-linear $\rho$ is usually expected in a system dominated by
electron-phonon scattering in the regime $T\gg\Theta_{D}$, where $\Theta_{D}$
is the Debye temperature; in the regime $T<\Theta_{D}$, the resistivity is
approximated by a low-order polynomial that is neither linear nor quadratic in
temperature. Transport measurements in the 122 compounds suggest
$\Theta_{D}\simeq 250$ K Ni _et al._ (2008); Tu _et al._ (2010). However,
the $T$-linear $1/\tau_{n}$ and $\rho_{n}$ spans a very broad temperature
range, from $T_{c}$ up to 300 K, which is inconsistent with the electron-
phonon dominated scattering process. Further clues on the $T$-linear
$1/\tau_{n}$ and $\rho_{n}$ may be revealed by an examination of the phase
diagram of the Ba1-xKxFe2As2 system and a comparison with similar behavior
found in other materials. The parent compound BaFe2As2 orders in an SDW state
below $T_{SDW}\approx 138$ K; $T_{SDW}$ is suppressed by K doping allowing
superconductivity to emerge. Here, $1/\tau_{n}\propto T$ is observed at the
doping where the SDW order is entirely suppressed, i.e. $T_{SDW}\rightarrow
0$. At this point, spin fluctuations are expected to be very strong, which has
been experimentally confirmed by NMR Ning _et al._ (2010); Nakai _et al._
(2010); Li _et al._ (2011). This brings us to models that attribute
$T$-linear $\rho$ to spin-fluctuation scattering Moriya and Ueda (2000);
Sachdev and Keimer (2011). The $T$-linear $\rho$ (or $1/\tau$) is also found
in cuprates such as Nd-doped La2-xSrxCuO4 Daou _et al._ (2009) and electron-
doped La2-xCexCuO4 Jin _et al._ (2011), organic superconductors (TMTSF)2PF6
Doiron-Leyraud _et al._ (2009a) and (TMTSF)2ClO4 Doiron-Leyraud _et al._
(2009b), as well as a number of heavy fermions such as CeCoIn5 Malinowski _et
al._ (2005); Tanatar _et al._ (2007) and YbRh2Si2 Trovarelli _et al._
(2000); Gegenwart _et al._ (2002). Studies on these materials have shown that
$T$-linear $\rho$, arising on the border of antiferromagnetic order, is caused
by spin-fluctuation scattering due to the proximity of an antiferromagnetic
QCP. Studies on cuprates and Bechgaard salts Taillefer (2010); Jin _et al._
(2011) further show that the strength (or coefficient) of the $T$-linear
$\rho$ scales with $T_{c}$ and disappears upon approaching the point where
$T_{c}\rightarrow 0$, suggestive of intimate relation between the $T$-linear
$\rho$ and superconductivity. Anomalous $T$-linear $\rho$ (or $1/\tau$) and
pairing in unconventional superconductors may share a common origin: spin
fluctuations. The hidden $T$-linear $1/\tau_{n}$ and $\rho_{n}$ in
Ba0.6K0.4Fe2As2 revealed by optical measurements may have the same origin as
those found in the cuprates, organic superconductors and heavy fermion metals,
because these materials share strikingly similar phase diagrams. Our
observations may also imply a possible QCP in the superconducting dome. The
existence of the QCP in iron-pnictides is supported by transport properties
Gooch _et al._ (2009); Shen _et al._ (2011); Maiwald _et al._ (2012), NMR
studies Ning _et al._ (2010); Nakai _et al._ (2010), de Haas-van Alphen
effect Walmsley _et al._ (2013), penetration depth measurement Hashimoto _et
al._ (2012) and first-principles calculations Xu _et al._ (2008).
Spin-fluctuation induced $T$-linear $\rho$ suggests an equivalent
$\omega$-linear $1/\tau(\omega)$ Moriya and Ueda (2000), obtainable through
the extended Drude model provided that interband contribution is negligible.
Low energy interband transitions are important in iron-pnictides, and their
contribution has to be subtracted to determine $1/\tau(\omega)$ for mobile
carriers Benfatto _et al._ (2011); Charnukha _et al._ (2011b). We calculated
$1/\tau(\omega)$ via the extended Drude model with the interband contribution
subtracted (supplementary) and found that by taking into account the interband
transitions, a large fraction of the frequency dependence in $1/\tau(\omega)$
is eliminated, which is consistent with Charnukha _et al._ ’s analysis on
Ba0.68K0.32Fe2As2 Charnukha _et al._ (2011b). There is no confident evidence
for the expected $\omega$-linear $1/\tau(\omega)$, since it could be masked by
the multiband character of the iron-pnictides.
In order to check if the $T$-linear $1/\tau$ is unique in Ba0.6K0.4Fe2As2 or
general in iron-pnictides, we applied the same analysis to BaFe2(As0.7P0.3)2
(supplementary). Interestingly, $T$-linear $1/\tau$ is also found for the
narrow Drude. This suggests that $T$-linear $1/\tau$ is not unique in
Ba0.6K0.4Fe2As2, but most likely, a general behavior in iron-pnictides at the
doping where SDW order is completely suppressed, i.e. $T_{SDW}\rightarrow 0$.
In summary, the detailed $T$ dependence of the normal state
$\sigma_{1}(\omega)$ and low-frequency spectral weight in Ba0.6K0.4Fe2As2 have
been examined. Two Drude components with different $1/\tau$’s yield an
excellent description of the low-frequency optical response, indicating the
existence of two groups of carriers with different quasiparticle lifetimes.
The broad Drude component produces an incoherent background conductivity with
no temperature dependence, while the narrow Drude component reveals a
$T$-linear $1/\tau_{n}$ and $\rho_{n}$. This fact explains the $T$-linear
$\rho$ behavior at low temperatures and the tendency to saturation at room
temperature observed in transport measurements in optimally hole-doped 122
compounds. An arctan($T$) dependence of the low-frequency spectral weight is
also a strong evidence for a $T$-linear $1/\tau$. Comparison with similar
behavior found in other materials suggests that the $T$-linear $1/\tau_{n}$
and $\rho_{n}$ in Ba0.6K0.4Fe2As2 may arise out of spin-fluctuation scattering
due to the possible existence of an antiferromagnetic QCP in the
superconducting dome.
We thank Hu Miao, Xiaoxiang Xi, Wei Ku, Cong Ren and Lei Shan for helpful
discussion. Work in Beijing was supported by the NSFC (No. 91121004 and No.
11104335) and the MSTC (973 Projects No. 2011CBA00107, No. 2012CB821400, No.
2012CB921302 and No. 2009CB929102). Work at BNL was supported by the U.S.
Department of Energy, Office of Basic Energy Sciences, Division of Materials
Sciences and Engineering under Contract No. DE-AC02-98CH10886. We acknowledge
the financial support from the Science and Technology Service of the French
Embassy in China.
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|
arxiv-papers
| 2013-09-11T16:33:34 |
2024-09-04T02:49:50.727722
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Y.M. Dai and B. Xu and B. Shen and H. Xiao and H.H. Wen and X.G. Qiu\n and C.C. Homes and R.P.S.M. Lobo",
"submitter": "Yaomin Dai",
"url": "https://arxiv.org/abs/1309.2881"
}
|
1309.2900
|
# Mining for Spatially-Near Communities in Geo-Located Social Networks
Joseph Hannigan
Dept. of Electrical Engineering
and Computer Science
U.S. Military Academy
West Point, NY
[email protected] &Guillermo Hernandez
Dept. of Electrical Engineering
and Computer Science
U.S. Military Academy
West Point, NY
Guillermo.Hernandez@
usma.edu &Richard M. Medina
Dept. of Geography
and GeoInformation Science
George Mason University
Fairfax, VA
[email protected] Patrick Roos
Dept. of Computer Science
University of Maryland
College Park, MD
[email protected] &Paulo Shakarian
Dept. of Electrical Engineering
and Computer Science
U.S. Military Academy
West Point, NY
[email protected]
###### Abstract
Current approaches to community detection in social networks often ignore the
spatial location of the nodes. In this paper, we look to extract spatially-
near communities in a social network. We introduce a new metric to measure the
quality of a community partition in a geolocated social networks called
“spatially-near modularity” a value that increases based on aspects of the
network structure but decreases based on the distance between nodes in the
communities. We then look to find an optimal partition with respect to this
measure - which should be an “ideal” community with respect to both social
ties and geographic location. Though an NP-hard problem, we introduce two
heuristic algorithms that attempt to maximize this measure and outperform non-
geographic community finding by an order of magnitude. Applications to
counter-terrorism are also discussed.
## Introduction
Community detection in social networks remains an important and active area of
research in the study of social network mining (?; ?; ?; ?; ?; ?; ?; ?; ?).
However, many real-world social networks also have a geographic context.
Social networks are tethered to geographic locations. People and their
relationships are tied to places. Even in the information age communications
are dependent on access. Though access can seem ubiquitous in many cases,
digital interaction cannot yet completely replace face-to-face contact,
especially for planned activities of spatiotemporal coincidence and the
transfer of tangible objects.
Primary considerations for research in many social science disciplines today
include characteristics of human activities and interactions in defined
spaces. The interactions can be between humans, or between humans and their
environments. These characteristics describe aspects of social complexity that
are necessary to understand when attempting to model open or closed social
systems. In studies of human security there is a new emphasis on
implementations of Activity Based Intelligence (ABI) to better understand
drivers toward specific actions and interactions, as well as to generate an
understanding of the system outside of targeted activities (?). Though the
concepts of ABI are new to many, its academic foundations of activity spaces,
social interaction, and spatio-temporal research are well established.
Attempts to identify sociogeographic based activity spaces, as demonstrated
here, are vital to the understanding of human behavior. Multi-spatial or
hybrid space (?) studies are much more valuable in this information age than
their single space counterparts. Multiple spaces are converging into hybrid
spaces as interactions in social systems become more complex.
In this paper, we look to develop a framework for deriving communities from
social network that is relevant not only with respect to network topology, but
also geography. The main geographic concept we use to relate nodes based on
space is “nearness.” On a general level, there exists a connection between
nearness and similarity. “Near” is a spatial concept, though not necessarily
geographically spatial. Social space nearness, or adjacency, typically
describes relationships between people or things that interact in some way.
Nearness based similarities need not be comprehensive. Single or few similar
traits can exist to maintain interaction; however, relatively more similar
traits between people can drive further or deeper interaction. Geographers and
sociologists have developed concepts that seek to explain the phenomenon of
nearness and similarity in their respective disciplines. In geography Tobler’s
First Law (TFL) describes this effect in physical space and homophily
describes it in social space (?; ?). Geographic and homophilic similarities
are inherently connected, as one of the greatest sources of homophily is
propinquity. Furthermore, interaction is driven by nearness and similarity.
The likelihood for interaction between people increases as distance decreases
between them. Community finding at the convergence of geographic and social
space nearness will lead to the identification of communities where place and
social traits drive interaction. Results of this method may be most meaningful
in studies of social systems that are greatly influenced by ethnicity and
culture among other geographically based factors. For example, communities
identified using geographic and social closeness may apply more to terrorist
and criminal networks than globally dispersed business networks.
Hence, our intuition is to find communities that are tightly-knit based on
network topology, but also spatially “near.” To do this, we create a new
measure of partition quality that we term “spatial nearness modularity” that
borrows concepts from network modularity (?) and $K$-means clustering (?; ?).
Hence, to find a high-quality set of communities with respect to this
geography and network connections, it stands to reason to search for an
optimal partition with respect to this measure. Unfortunately, we are able to
show that doing so is NP-hard based on the results of (?). To address this
issue of intractability, we introduce two heuristics and we then
experimentally evaluate them, where we find that our approach provides an
order-of-magnitude improvement in spatially-near modularity over non-
geographic approaches. This is followed by a description of how this technique
could apply to counter-terrorism and a discussion of related work.
## Technical Preliminaries
We assume the existence of a undirected graph $G=(V,E)$ where set $V$ are
vertices and $E$ are edges among them. As the graph in undirected, $(i,j)\in
E$ implies $(j,i)\in E$. We shall use $n,m$ to represent the sizes of $V,E$
respectively. Each edge $(i,j)$ will be associated with a positive real weight
denoted by $w_{ij}$ (if there is no edge between $i$ and $j$, $w_{ij}=0$). For
a given node $i\in V$, we shall use the symbol $\eta_{i}$ to represent the set
$\\{j\in V|\exists(i,j)\in E\\}$ and $k_{i}$ is the size of this set. We shall
also assume a distance function $d:V\times V\rightarrow\Re$ that meets the
normal distance axioms: $d(i,i)=0$, $d(i,j)=d(j,i)$, and $d(i,j)\leq
d(i,j^{\prime})+d(j,j^{\prime})$. For ease of notation, we shall use $d_{ij}$
instead of $d(i,j)$. In this paper we will often use the notation
$C=c_{1},\ldots,c_{x}$ to denote a partition of $V$. Hence, $\cup_{c_{i}\in
C}c_{i}=V$ and for all $c_{i},c_{j}\in C$ $c_{i}\cap c_{j}=\emptyset$. We
define the modularity of a partition ($M(C)$) in accordance with the
definition introduced by (?) as follows:
NG-Modularity. (?) Given a social network $G=(V,E)$ and partition $C$ the
Newman-Girvan (NG) modularity is defined as follows:
$\displaystyle M_{NG}(C)=\frac{1}{2m}\sum_{c\in C}\sum_{i,j\in
c}w_{ij}-\frac{k_{i}k_{j}}{2m}.$ (1)
The modularity of a network partition measures the quality of its partition
structure as the density of edges within partitions compared to the density of
edges between partitions. The former is ideally very high compared to the
latter. Modularity will give a number in $[-1,1]$, a higher value meaning
better quality partition. Previous work, such as (?; ?), has focused on
finding a partition that optimizes this quantity. However, modularity
maximization only considers network topology and does not make any effort to
group individuals that are geographically close to each other. An alternative
is to find a partition of $K$ clusters of nodes that minimizes the sum-of-
squares distance to the center of each cluster. This is known as $K$-means
clustering (?; ?). $K$-means clustering algorithms attempt to find a partition
of points on a plane into $K$ clusters, such that the following quantity is
minimized (here $x_{c}$ is the centroid of the points in cluster $c$):
$\displaystyle\sum_{c\in C}\mathbf{agg}_{i\in c}d(i,x_{c})^{2}.$ (2)
In the above definition, $\mathbf{agg}$ is some aggregate function. Common
aggregates used here are $\max$ and $\sum$. For the purpose of this paper, as
modularity is maximized, we wish to minimize some aggregate of the distances
to the center of each cluster. Thus, one potential quantity that could be
optimized is the following:
$\displaystyle\frac{1}{2m}\big{(}\frac{\sum_{c\in C}\sum_{i,j\in
c}w_{ij}-\frac{k_{i}k_{j}}{2m}}{1+\sum_{c\in C}\mathbf{agg}_{i\in
c}d(i,x_{c})^{2}}\big{)}.$ (3)
Note that the additive $1$ in the denominator is to avoid division by zero and
to ensure that the result will be within the range $[-1,1]$. The above
optimization function has the useful property that we can embed both
modularity maximization and $K$-means clustering - the first by placing all
nodes in the same location, the second by ignoring edges among any nodes in
the network and restricting the number of clusters to be exactly $K$. However,
one aspect the above definition misses is that it cannot measure the quality
of an individual community. Hence, we introduce an alternative definition
below that we term “spatially-near (SN) modularity.”
SN-Modularity. Given a social network $G=(V,E)$, partition $C$, and scaling
parameter $\sigma\in\Re^{+}$, the spatially-near (SN) modularity is defined as
follows:
$\displaystyle M_{SN}(C,\sigma)=\frac{1}{2m}\sum_{c\in C}\frac{\sum_{i,j\in
c}w_{ij}-\frac{k_{i}k_{j}}{2m}}{1+\mathbf{agg}_{i\in
c}\left(\frac{d(i,x_{c})}{\sigma}\right)^{2}}.$ (4)
So, for a given community, we can measure its quality with the following:
$\displaystyle\frac{1}{2m}\frac{\sum_{i,j\in
c}w_{ij}-\frac{k_{i}k_{j}}{2m}}{1+\mathbf{agg}_{i\in
c}\left(\frac{d(i,x_{c})}{\sigma}\right)^{2}}.$ (5)
We also note that as $\sigma$ increases, distance is de-emphasized. This
parameter would be specified based on the relative importance of distance to
to network structure as well as the unit of measurement used for distance.
Practically, a user could potentially provide this parameter in many different
ways. Simple methods would include setting $\sigma$ to $1$, the average
distance among all pairs of nodes, or the average distance among all pairs of
nodes that have an edge between them. Alternatively, this parameter could also
learned from historical data, if such a corpus is available. Another approach
is for the user to explore various parameter settings. In this work, we leave
advanced methods for determining $\sigma$ to future work and conduct
experiments with multiple settings for this parameter. However, we note that
for particularly large values of $\sigma$, SN-modularity becomes equivalent to
NG-modularity. It is easy to show the following property:
$\displaystyle\mathit{lim}_{\sigma\rightarrow\infty}M_{SN}(C,\sigma)=M_{NG}(C).$
(6)
However, maximizing $M_{SN}(C,\sigma)$ remains NP-hard. Hence, in this paper
we introduce two heuristic algorithms to find a partition $C$ where
$M_{SN}(C)$ is near-optimal.
Theorem. For a given social network $G=(V,E)$ and scaling factor $\sigma$,
identifying a partition $C$ s.t. $M_{SN}(C,\sigma)$ is maximized is NP-hard.
Proof. We can embed an instance of finding a partition $C$ that maximizes
$M_{NG}(C)$ into the problem from the statement by creating a distance
function $d$ where $\forall i,j,d(i,j)=0$ and setting $\sigma$ to an arbitrary
value. Hence, any algorithm that maximizes $M_{SN}$ using this construction
also maximizes $M_{NG}$. Since finding a partition that maximizes $M_{NG}$ is
NP-hard by the results of (?), the statement of the theorem follows.
$\blacksquare$
## Algorithms
In the previous section, we found that identifying a spatially-near partition
is an NP-hard problem. Hence, in this section, we propose two heuristic
approaches to deal with this intractability. We later describe our evaluation
of these approaches. In our first heuristic, which we call “Louvain-SN”, we
employ the modification of the Louvain algorithm of Blondel et al. (?), only
instead of using it to maximize NG-modularity, we use it to maximize SN-
modularity (the Louvain algorithm was designed to find a near-optimal parition
w.r.t. NG-modularity). Our second algorithm, the SNIC (Spatially Near
Iterative Constraining) algorithm, relies on multiple calls to the Louvain-SN
algorithm - but each with a limit on the aggregate distance permitted in a
community.
### The Louvain-SN Algorithm
The original Louvain algorithm of (?) is an iterated, hierarchical process in
which two phases are applied repeatedly until maximal modularity is reached:
During the first phase, each node $v_{i}\in V$ of the given social network is
assigned to a community $c$, creating an initial partition. In (?), the
singleton partition was used - which we use in this work as well. Then, for
each $v_{i}\in V$, the gains in modularity that would result from moving
$v_{i}$ to the community of each of its neighbors $v_{j}\in\eta_{i}$ are
calculated, and $v_{i}$ is removed and placed into the community for which the
maximum improvement in modularity is achieved (unless no positive gain in
modularity is possible). This sub-process is repeated sequentially for each
$v_{i}\in V$ until no individual move will result in a gain in modularity,
marking the end of the first phase and giving a partition $C$. During the
second phase, a new network is built by using each $c_{i}\in C$ as a node in
the new network, call these nodes meta-nodes. Weights on the edges between any
two meta-nodes in the new network are assigned to be the sum of the weights of
the edges between nodes in the two communities corresponding to the meta-
nodes. Here, self-loops are created for each meta-node in the new network from
the links between nodes of the community corresponding to that meta-node.
After this phase is complete, the two phases are reapplied iteratively until
there are no more changes.
The efficiency of the Louvain algorithm relies on an easy re-calculation of
modularity in the first phase of the algorithm. When computing gains in
modularity in phase one of the algorithm, removing any node $v_{i}$, the
overall increase in modularity if it is placed into community $c$ is
proportional to:
$\displaystyle k_{i,in}-\sum_{j\in c}\frac{k_{i}k_{j}}{2m}$ (7)
In our modification for optimizing SN-modularity, we can retain some of this
efficiency by retaining the previous denominator and numerator of Equation 5
(multiplied by $2m$) for each community. By retaining these values along with
the value of Equation 7, computing the increase or decrease in modularity for
a community can be performed quickly (though this ultimately depends on how
the aggregate function $\mathbf{agg}$ and the centers of the communities
$x_{c}$ are computed). Additionally, in the creation of the meta-nodes, we use
the centers from the previous step as their location. Additionally, we also
found that we obtained improvement in performance by allowing a removed node
to be moved back to a community containing just itself, as unlike in the
maximization of standard modularity, isolating a node in this fashion could
potentially increase the overall modularity due to the denominator of
Expression 4.
### The SNIC Heuristic
Next, we introduce the “Spatially Near, Iterative Constraining” (SNIC)
Heuristic. This idea was created as the result from a pilot experiment where
we noticed that constraining a node to join only communities where it was
geographically “near” to all the members would sometimes improve the resulting
quantity of $M_{SN}$. The question is how does one determine where to set this
distance constraint. In our experiments we ran our modified Louvain approach
iteratively, returning only the maximum distance between two points in the
community upon each iteration. This distance is then set as the distance
constraint for the next iteration. Once the distance constraint reaches zero
(or a maximum number of iterations is reached), the algorithm then returns the
partition found which is associated with the greatest value for Expression 4.
## Experimental Results
For our experiments, we used information extracted from the Brightkite
location-based online social networking sites (?).
We built our implementation in Python 2.6 on top of the NetworkX
library111http://networkx.github.com/ leveraging code from Thomas Aynaud’s
implementation of the Louvain
algorithm222http://perso.crans.org/aynaud/communities/. Our implementation
took approximately $1000$ lines of code. The experiments were run on a
computer equipped with an Intel Core i7 Processor operating at 2.67 GHz (one
core utilized) running Microsoft Windows 7 and equipped with 4.0 GB of
physical memory. All statistics presented in this section were calculated
using SPSS 19. We use our heuristics to find partitions based on Expression 4
where $\mathbf{agg}=\max$.
In our first set of tests, we iteratively selected nodes and their neighbors
from the Brightkite network dataset provided by the authors of (?) to produce
10 samples of 1000 nodes. To generate the samples, each sample begins with a
randomly selected node from the network. The selected node and all of its
connected nodes are then included in the next iteration, in which a new random
node is chosen. This continues until 1000 nodes are reached for each sample.
The minimum edge count for all samples processed is 1729, while the maximum is
2282. The average number of edges is 1929.
In our trials, we varied the $\sigma$ parameter with the values
$\\{300,500,1000,2000,3000,4000,5000\\}$. For each dataset and each value of
$\sigma$, we compare the SN-modularity returned by three approaches: the
Louvain algorithm (does not consider geospatial information), the Louvain-SN
algorithm (the modified version of the Louvain algorithm for SN-modularity
optimization), and the SNIC algorithm (an iterated version of the Louvain-SN
that selects the best result based on updating the distance constraint).
The SNIC algorithm returned a partition with greater average SN-modularity for
each value of $\sigma$ than the partitions returned by the Louvain and
Louvain-SN algorithms (see Figure 1). In general, the SNIC algorithm
consistently outperformed the Louvain algorithm in terms of SN-modularity -
producing a partition of greater SN-modularity on all trials. The Louvain-SN
outperformed the standard Louvain in all but $11$ (of $70$) trials, though (as
we discuss later in this section) this improvement is likely not statistically
significant, unlike the SNIC heuristic.333There were $11$ such trials out of
the $70$ trials where the Louvain outperformed the Louvain-SN. Of the cases
where there was decreased quality over standard Louvain, the maximum decrease
in quality was $26.52\%$ and the average decrease was $15.10\%$. The SNIC
algorithm outperformed the Louvain algorithm on all trials.
Figure 1: $\sigma$ (in kilometers) vs. (average) SN-modularity for the
partitions returned by the Louvain, Louvain-SN, and SNIC algorithms.
To determine significant difference in SN-modularity of the three approaches
on the Brightkite dataset, analysis of variance (ANOVA) tests were used.
Difference in SN-modularity for the three approaches was confirmed with a
p-value of 0.006. A Tukey’s Honest Significant Difference (HSD) test was also
used to determine pairwise differences between the approaches. No significant
differences were found between the Louvain and Louvain-SN algorithms; however,
the SNIC algorithm was found to be different than both the Louvain (at p =
0.010) and the Louvain-SN (at p = 0.020). Additionally, the differences for
runtimes of the three approaches were found to be different with a p-value of
0.000 (see Figure 2). As with the difference in SN-modularity, runtime
differences exist between the Louvain and the SNIC algorithms (at p = 0.000)
and between the Louvain-SN and the SNIC algorithms (at p = 0.000). These
results are also provided through use of Tukey’s HSD. Differences in SN-
modularity and runtimes for the three approaches can be seen in Figures 1 and
2, respectively. Further, we also note that although the SNIC algorithm has
significantly greater runtime than the Louvain and Louvain-SN, it still appear
to scale linearly with the number of nodes in the network ($R^{2}=0.992$).
Hence, it may still be a viable solution for very large networks. We are
currently studying the scalability of this algorithm.
Figure 2: Network size (by nodes) vs. runtime for the partitions returned by
the Louvain, Louvain-SN, and SNIC algorithms.
Figure 3 shows the increase in quality of community finding (SN-modularity)
over iterations of the SNIC algorithm. Recall that the SNIC algorithm
decreases the distance constraint at each iteration. As the geographic
constraint decreases, such that community proximity becomes more important,
the quality of community (number of connections within vs outside) increases.
Here we introduce an axiom - that as the geographic space of interaction for a
social network shrinks, it is more likely that those left within the community
are more connected. Spatial outliers, which are also social outliers can be
conceptualized as weak links (?) and are removed through community proximity
limiting iterations. Through 100 iterations, the quality of community
increases and in most social networks this value may continue to increase
given high enough spatial resolution data. In other words, humans form
communities and interact mostly with those they are geographically near, such
that the strongest communities will be those shared within small geographic
proximities. However, we note that more iterations of the SNIC algorithm will
not result in singleton communities, as that is the initial partition
considered by the algorithm. Also note that the improvement in SN-modularity
as a function of number of iterations of the SNIC heuristic is also likely
dependent on the parameter $\sigma$.
Figure 3: SN-modularity vs. number of iterations for the partitions returned
by the SNIC algorithms for the Britekite network data.
As with the example above in Figure 3, Figure 4 shows that a stronger
influence of geographic distance on community finding leads to greater quality
communities based on the SN-modularity measure. Recall that $\sigma$ is the
scaling parameter in SN-modularity. Decreasing the scaling parameter, in turn
strengthening the geographic influence on the equation, leads to an increase
in the quality of communities identified by the SNIC algorithm. Increasing the
$\sigma$ value will result in an asymptotic trend for SN-modularity toward
that expected from the non-spatial Louvain algorithm. This trend is shown for
the Brightkite network in percent increase in SN-modularity over the Louvain
algorithm for (A) the Louvain-SN and (B) the SNIC algorithms.
Figure 4: $\sigma$ (in kilometers) vs. percent improvement in SN-modularity
for the partition returned by the Louvain-SN (panel A) and SNIC (panel B)
algorithms. Not depicted in panel A (Louvain-SN) are results where the
Louvain-SN algorithm produced lower-quality results than the standard Louvain
(due to the log-scale, see text for further details).3 Note that for the SNIC
algorithm (panel B) outperformed the standard Louvain on all trials.
The difference in Brightkite communities identified by the Louvain and SNIC
algorithms is clear in Figure 5. There is a quantitative difference as
suggested by the SN-modularity metric results, but also a very qualitative
difference in which the communities identified by the SNIC algorithm are much
more spatially constrained. The bottom half of Figure 5 represents the SNIC
algorithm results with $\sigma$ = 1. In today’s information age where global
networks are common, methods to identify geograhically unconstrained
communities, as well as those methods that identify their geographically
constrained counterparts are both equally valuable. Implications for strength
of ties, activity and operations spaces, and interactions are different when
considering geographic network characteristics.
Figure 5: Top: Brightkite Communities identified using the Louvain algorithm,
Bottom: Communities identified using the SNIC algorithm.
Additionally, we also studied the NG-modularity of the partition returned by
the SNIC algorithm. We found that although the SNIC algorithm was not designed
to maximize NG-modularity, it still provided a positive value - which
indicates that there is a greater density of edges within the communities as
opposed to between communities (Figure 6). We also found that the solution
returned by the NG-modularity of the partition returned by the SNIC algorithm
seems to approach the NG-modularity of the solution of the Louvain algorithm
as $\sigma$ increases. Although this is not guaranteed theoretically, it
should be expected based on the relationship between NG-modularity and SN-
modularity shown in Equation 6.
Figure 6: NG-modularity of the partition returned by the SNIC algorithm ($10$
iterations) as a function of $\sigma$.
## Applications
There are many fitting applications for algorithms that detect sociogeographic
communities. In general, any network that requires or benefits from geographic
propinquity can serve as a test case for the SNIC algorithm. For example, the
diffusion of a disease through a social network requires geographic closeness,
or face-to-face interaction, between people. While much of the diffusion may
not be social network based, but solely spatial, those that have stronger
social ties and in turn interact more in geographic space are more likely to
contract or spread a disease. This phenomenon exists at various levels of
physical interaction for contagious diseases. This model of diffusion works
with the spread of any biologically contagious, material, or even ideological
transfer that requires coincidence in space and time. The following example
shows the value of the SNIC algorithm for sociogeographic analysis on a
transnational terrorist network.
Figure 7 illustrates the difference between community finding results using
both the Louvain (non-spatial) and SNIC algorithms on a transnational
terrorist network dataset. The dataset used for this research is
representative of a global Islamist terrorist network from the late 1970s to
approximately 2010 (see (?) for a full description and discussion of the
dataset). The SNIC algorithm application shown here uses $\sigma$ = 1600. The
transnational Islamist terrorist network is a cellular based, decentralized
structure and heavily dependent on relative location and proximity (?).
Because this is the case, identifying sociogeographic communities requires
only a small spatial component additional to the social component.
Applications of the SNIC algorithm on other network structures may require
more spatial influence to identify sociogeographic communities. For example,
the Brightkite application shown in Figure 5 uses $\sigma$ = 1. The terrorist
network is much smaller with 358 nodes and 660 edges, and is much more
geographically based for operational necessity.
As stated previously, research results that identify social closeness vs.
those that identify sociogeographic closeness are quantitatively and
qualitatively different for many social networks. The top graphic in Figure 7
shows the modularity results using the Louvain algorithm, which highlight the
transnational network connections. Many of the Islamist terrorist network
cells have foundations or affiliates in Europe and the Middle East. While the
strength of social communities can be equal over long distances, especially if
network connections were made at some point coincident in space and time, it
is beneficial to isolate communities in geographic space for some
applications. In this case, the SNIC algorithm successfully identifies
operational communities (A) the 9/11 cell planning and preparing for the
attack in Southern California and Arizona, (B) a father and son diad working
with, specifically financing, al-Qaeda in Canada, (C) a sociogeographic
community of al-Qaeda tied members in Montreal, Canada, some of which were
plotting to attack Los Angeles International Airport in 1999, (D) two al-Qaeda
linked cells in Boston, MA with members in the Boston sleeper cell and
plotting a large scale bombing attack in Jordan at multiple sites, (E) the al-
Qaeda based cell operating in New York responsible for the first World Trade
Center Attack in 1993, and (F) communities of 9/11 hijackers operating in
Florida and other eastern US states. The SNIC algorithm can be additionally
adjusted to further separate cellular communities based on geography (by
varying $\sigma$ and the number of iterations of the algorithm).
In systems such as this terrorist network, connected individuals that are
close in geographic space, but not as close socially, can be more important to
identify when attempting to counter operations. For example, identification of
weaker but closer social links, such as those providing materials to a
terrorist cell can be used as valuable intelligence to understand and
dismantle terrorist operations in local to regional settings. Knowledge of
international connections is important for understanding the global terrorist
system, and cells in decentralized networks often maintain communications over
long distances. However, many of these cells can operate independently, though
in most cases they will need proximal resources. These local system
interactions can be detected through use of the SNIC algorithm.
Figure 7: Top: Terrorist communities identified using the Louvain algorithm,
Bottom: Terrorist communities identified using the SNIC algorithm.
## Related Work
Modularity maximization for community finding was first introduced in (?). In
(?), the Louvain algorithm is introduced, which can scale to very large
networks and is shown to provide partitions that nearly maximize modularity.
We leverage a modification of the Louvain algorithm in this paper. Finding
geographically disperse communities in a social network has also been
previously studied (?; ?; ?; ?). Our approach in this paper differs in that we
desire to find communities where the nodes are spatially-near and not distant.
In addition to the aforementioned approaches, community detection in networks
has also been explored in other manners that have potential to be applicable
to the geospatial case - though to our knowledge no such application has been
presented in the literature. See (?; ?) for examples.
There also exist many approaches for community detection in networks not based
on modularity maximization. Examples use label propagation (?), random walks
(?), or bottom-up voting approaches (?). See (?) for comprehensive surveys.
These do not consider spatial interactions - leveraging these approaches in a
geospatial context is an important possibility for future work.
Geospatial networks have been explored with respect to problems other than
community finding such as link-prediction (?) and identifying user location
(?). There have also been several empricial studies on social networks with a
spatial component such as (?; ?; ?). More domain-specific empirical studies
related to this work are also prevalent in the literature. Pertinent to our
application are studies on terrorist networks (?) and criminal co-offender
networks (?).
## Conclusion
In this work, we introduced spatially-near modularity - a measure of the
quality of a geographically-near partition in a social network. Though finding
an optimal partition with respect to this measure is NP-hard, we were able to
obtain quality partitions with two heuristic algorithms that we introduced in
this paper and tested on real-world datasets. We have also discussed various
ways in which our algorithms can be applied to gain useful knowledge in
counter-terrorism applications. Our immediate concern for future work is
exploring the scalability of this approach ($10^{6}$ nodes and greater).
Additionally, we are also pursuing temporal dynamics of such communities and
the differences between the communities formed based on the current state of
the nodes (i.e. “work” vs “home”). In our more practical research, we are also
working to integrate the generation of geographically-near partitions into our
Organizational, Relationship, and Contact Analyzer (ORCA) software (?) that we
are currently fielding to several American law-enforcement agencies.
Acknowledgements. This work was supported by ARO (project 2GDATXR042 and grant
W911NF-08-1-0144), AFOSR (grant FA9550-12-1-0021) and the Office of the
Secretary of Defense. The opinions in this paper are those of the authors and
do not necessarily reflect the opinions of the funders, the U.S. Military
Academy, or the U.S. Army.
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|
arxiv-papers
| 2013-09-11T17:34:43 |
2024-09-04T02:49:50.735751
|
{
"license": "Public Domain",
"authors": "Joseph Hannigan, Guillermo Hernandez, Richard M. Medina, Patrcik Roos,\n Paulo Shakarian",
"submitter": "Paulo Shakarian",
"url": "https://arxiv.org/abs/1309.2900"
}
|
1309.2904
|
# A System-Theoretic Clean Slate Approach to Provably Secure Ad Hoc Wireless
Networking
Jonathan Ponniah*, Yih-Chun Hu*, and P. R. Kumar *CSL & ECE, Univ. of
Illinois, 1308 West Main St., Urbana, IL 61801. Email:
{ponniah1,yihchun}@illinois.edu. Tel: 217-333-4220. **Corresponding author:
ECE, Texas A&M University, 3259 TAMU, College Station, TX 77843-3259. Email:
[email protected]. Tel: 979-862-3376. This paper is partially based on work
supported by NSF under Contract Nos. CNS-1302182, CCF-0939370 and CNS-1232602,
AFOSR under Contract No. FA-9550-13-1-0008, and USARO under Contract No.
W911NF-08-1-0238.
###### Abstract
Traditionally, wireless network protocols have been designed for performance.
Subsequently, as attacks have been identified, patches have been developed.
This has resulted in an “arms race” development process of discovering
vulnerabilities and then patching them. The fundamental difficulty with this
approach is that other vulnerabilities may still exist. No provable security
or performance guarantees can ever be provided.
We develop a system-theoretic approach to security that provides a complete
protocol suite with provable guarantees, as well as proof of min-max
optimality with respect to any given utility function of source-destination
rates. Our approach is based on a model capturing the essential features of an
ad-hoc wireless network that has been infiltrated with hostile nodes. We
consider any collection of nodes, some good and some bad, possessing specified
capabilities vis-a-vis cryptography, wireless communication and clocks. The
good nodes do not know the bad nodes. The bad nodes can collaborate perfectly,
and are capable of any disruptive acts ranging from simply jamming to non-
cooperation with the protocols in any manner they please.
The protocol suite caters to the complete life-cycle, all the way from birth
of nodes, through all phases of ad hoc network formation, leading to an
optimized network carrying data reliably. It provably achieves the min-max of
the utility function, where the max is over all protocol suites published and
followed by the good nodes, while the min is over all Byzantine behaviors of
the bad nodes. Under the protocol suite, the bad nodes do not benefit from any
actions other than jamming or cooperating.
This approach supersedes much previous work that deals with several types of
attacks including wormhole, rushing, partial deafness, routing loops, routing
black holes, routing gray holes, and network partition attacks.
###### Index Terms:
Ad hoc wireless networks, security.
## I Introduction
Our focus is on the problem of security of ad-hoc, multi-hop, wireless
networks. The wireless nodes in these types of networks need to determine when
to transmit packets and at what power levels, discover routes from sources to
destinations, and ensure overall end-to-end reliability, all without any
centralized controller guiding the process. This requires a suite consisting
of multiple protocols.
Several candidate have been proposed. Medium access control protocols include
IEEE 802.11 [10] and MACAW [2], power control protocols include COMPOW [12]
and PCMA [15], routing protocols include DSDV [17], AODV [16], DSR [11], and
OLSR [5], and transport protocols include TCP [20] and variations for ad hoc
networks [13, 6, 3, 21].
All the above protocols are designed on the assumption that all nodes are
“good,” and will conform to the protocol. Some nodes can however be malicious,
deliberately intent on disrupting the network, a vulnerability especially
acute since the very purpose of ad hoc networks is to allow any node to join a
network. For wireless networks used in safety-critical applications, e.g.,
vehicular networks, vulnerabilities can be dangerous. Moreover, many wireless
networking protocols have been based on wireline protocols, with possible
susceptibilities to novel over the air attacks.
The assumption of benignness, implicit or explicit, has been the traditional
starting point of protocol development. Systems have been first designed to
provide high performance. Subsequently, as vulnerabilities have been
discovered, they have been patched on a case by case basis. For example, the
“wormhole” attack was discovered in [7], where an attacker sets up a false
link between two nodes. It is countered by a fix using temporal and
geographical packet leashes [7, 19]. The “rushing” attack against DSR was
discovered in [8], in which attackers manipulate the network topology. This is
countered by a fix using network discovery chains. The “partial deafness”
attack against 802.11 was discovered in [4], in which an attacker artificially
reduces its link quality to draw more network resources. It is countered by a
fix using queue regulation at the access point. Other attacks against DSR are
the routing loop attack in which an attacker generates forged routing packets
causing data packets to cycle endlessly; the routing black hole attack in
which an attacker simply drops all packets it receives; and the network
partition attack in which an attacker injects forged routing packets to
prevent one set of nodes from reaching another. These attacks are all
countered in the Ariadne protocol [9] by the joint use of routing chains,
encryption, and packet leashes. Some protocols such as Watchdog and Pathrater
[14] try to pre-empt attacks by maintaining a blacklist that tracks malicious
behavior, but this backfires if an attacker maligns a good node, causing other
good nodes to add that node to their blacklists. These attacks are not
targeted at violating privacy of communications between nodes, which can be
avoided simply by encryption. Rather, they are generally Denial of Service
attacks (DoS), which usually take advantage of algorithms that assume the
participating users are good or cooperative.
The basic problem with this arms race approach of hardening algorithms
initially designed for good performance is that one never knows what other
vulnerabilities or attacks exist. Thus no guarantees can be provided about the
security of the protocols at any stage of the arms race process.
Our goal in this paper is to propose an alternate clean slate system-theoretic
approach to security that provides provable performance guarantees. We pursue
a model-based approach, comprising a physical model of node capabilities,
clocks, cryptography, and wireless communication. It is an initial attempt to
holistically model the entire dynamics of an ad-hoc wireless network that has
been infiltrated with hostile nodes.
Our goal is to design a protocol suite for the complete life-cycle of the
wireless system, all the way from the very birth of the nodes, and continuing
through all phases of the network formation process, to a long-term operation
where the network is carrying data reliably from sources to their
destinations. The good nodes don’t know who the bad nodes are, and are
required to follow the published protocol suite. Throughput all phases, the
bad nodes can perfectly collaborate and incessantly indulge in any disruptive
behavior to make the network formation and operation dysfunctional. They could
just “jam,” or engage in more intricate behavior such as not relay a packet,
advertise a wrong hop count, advertise a wrong logical topology, cause packet
collisions, disrupt attempts at cooperative scheduling, drop an ACK, refuse to
acknowledge a neighbor’s handshake, or behave inconsistently.
We design a protocol suite that is provably secure against all such attacks by
the malicious nodes. Not only that, it guarantees min-max optimal performance.
The performance is described by a given utility function, which the good nodes
wish to maximize by publishing a complete protocol suite and conforming to it.
The bad nodes on the other hand aim to minimize this utility by indulging in
all manner of “Byzantine” behavior described above not conforming to the
protocol.
This leads to a zero-sum game. Since the good nodes first announce the
protocol, the best value of the utility function that the good nodes can hope
to attain is its max-min, where the maximization is over all protocol suites,
and the minimization is over all Byzantine behaviors of the bad nodes. We will
prove that the protocol suite designed attains this max-min to within any
$\epsilon>0$. Moreover, we establish three even stronger results.
First, this game actually has a saddle point, i.e., the protocol suite attains
the min-max (to within any $\epsilon>0$). (Generally, min-max results in a
higher utility than max-min, since the bad nodes have to first disclose their
tactics).
Second, the bad nodes can do no better than just jamming or conforming to the
published protocol suite on each “concurrent transmission vector,” a
generalization of the notion of an “independent set” of nodes that can
simultaneously transmit. They do not benefit from more elaborate Byzantine
antics.
Third, the protocol optimally exploits any non-hostile behavior of the bad
nodes. If they behave suboptimally, i.e., are not as hostile as they could be,
then it will take advantage. This is a desirable feature since while one wants
to design protocols that are guaranteeably secure in the worst case, one would
want them to exploit any benignness in the environment.
Some important qualifications need to be noted. First, the results are valid
only for the postulated model of the network. Future research may identify
technological capabilities outside the model that can attack the protocol
suite. Such discoveries will, one hopes, lead to the development of more
general models and procotols provably secure in them. The research enterprise
will thereby be elevated to a higher level; instead of reacting to each
proposed protocol one reacts to each proposed model, with provable guarantees
provided at each step. Section VII provides some such directions for model
generalization.
Second, though not merely asymptotic, the optimality is over a large time
period, and the overhead of transient phases of the protocol may be high.
However, there is much scope for optimizing protocol overhead while preserving
security.
Third, how should one view the proposed protocol suite? The answer is layered.
At a minimum, it can be regarded as a constructive existence proof that one
can indeed provide optimal performance while guaranteeing security, with the
identified model class only serving as an exemplar of conditions under which
this can be done. To a more receptive reader, the designed protocol suite is
suggestive of of how one can do so. The architectural decomposition into
several phases could perhaps be kept in mind by future protocol designers.
At any rate, one hopes that this approach will trigger several critical
reactions among a skeptical readership, and lead to follow up work that
designs protocols with guaranteed security and performance for more general
model classes.
Section II describes the model, Section III the main results, Section IV an
outline of the approach, Section V the protocol suite, and Section VI proves
feasibility and optimality.
## II The Model
The model of an ad-hoc wireless network infiltrated by hostile nodes can be
organized into four categories: the nodal model (N), communication model (CO),
clock behavior (CL), and cryptographic capabilities (CR).
Nodal model: (N1) There are $n$ nodes, some good and some bad. Let $G$ denote
the set of good nodes, and its complement $B$ the set of bad nodes. (N2) The
good nodes do not know who the bad nodes are a priori. (N3) The bad nodes are
able to fully coordinate their actions, and are fully aware of their
collective states (equivalent to unlimited bandwidth between them). (N4) The
good nodes are all initially powered off, and they all turn on within $U_{0}$
time units of the first good node that turns on.
Communication model: (CO1) Each node $i$ can choose from among a finite set of
transmission/reception modes $M_{i}$ at each time. Each mode corresponds, if
transmitting, to a joint choice of power level, modulation scheme and encoding
scheme for each other intended receiver node, or to just listening and not
transmitting, or even to “jamming,” which simply consists of using its power
output to emit noise. (CO2) The good nodes are half-duplex, i.e., cannot
transmit and receive simultaneously. (CO3) We call
$c=(c_{1},c_{2},\ldots,c_{n})$ denoting the mode choices of all the nodes made
at a certain time, as a “concurrent transmission vector” (CTV). (It is more
general than an independent set that is sometimes used to model wireless
networks). We will denote by $c_{G}=(c_{i}:i\in G)$ and $c_{B}=(c_{i}:i\in B)$
the vectors of choices of modes made by the good and bad nodes respectively,
with each $c_{i}\in{\cal{M}}_{i}$, and let ${\cal{C}}_{G}$ and ${\cal{C}}_{B}$
denote the sets of all such choices. We will denote by
${\cal{C}}:={\cal{C}}_{G}\times{\cal{C}}_{B}$, the set of all CTVs. (CO4) Each
$c$ results in a “link-rate vector” $r(c)$ of dimension $n(n-1)$. Its $ij$-th
component, $r_{ij}(c)$, is the data rate at which bits can be sent from node
$i$ to node $j$ at that time. Due to the shared nature of the wireless medium,
the rate depends on the transmission mode choices made by all the other nodes,
as well as the geographic locations of the nodes, the propagation path loss,
the ambient noise, and all other physical characteristics affecting data rate.
A component $r_{ij}(c)$ may be zero, for example if the SINR at $j$ is below a
threshold value for decoding, or if node $i$ is not transmitting to node $j$.
(CO5) If a certain rate vector is achievable then lower rates are also
achievable. To state this, let $\Lambda:=\\{r_{ij}(c):i\neq
j,c\in{\cal{C}}\\}$ denote the finite set of all possible rates than can be
achieved. We suppose that for every $c$, and $r^{\prime}\leq r(c)$ (understood
component wise) with all elements in $\Lambda$, there is a choice
$c^{\prime}\in{\cal{C}}$ such that $r(c^{\prime})=r^{\prime}$. This assumption
is not strictly necessary, but it helps to simplify the statement that bad
nodes can claim to receive only at low rates. (CO6) In the case of a bad node
$j$, the rate $r_{ij}(c)$ may be the result of some other bad node being able
to decode the packet from $i$ at that rate, and then passing on that packet to
$j$, since bad nodes can collaborate perfectly. In the case of a bad node $i$,
the rate $r_{ij}(c)$ may be the result of some other bad node being able to
transmit the packet successfully to $j$ at that rate, pretending to be $i$.
Meanwhile, in either case, the bad node may be jamming. Thus a bad node can
both jam and appear to be cooperating, whether transmitting or receiving, at
the same time. (CO7) The bad nodes can claim to have received transmissions
from each other at any of the rates in the finite set $\Lambda$, as they
please. To state this, for $c=(c_{G},c_{B})$, we will partition the resulting
link-rate vector as $r(c)=(r_{GG}(c),r_{GB}(c),r_{BG}(c),r_{BB}(c))$, where
$r_{BG}$ denotes the link-rates from the bad nodes to the good nodes, etc. We
suppose that for every $c=(c_{G},c_{B})$ and every $r^{\prime}$ with all
elements in $\Lambda$, there is a $c^{\prime}_{B}\in{\cal{C}}_{B}$ such that
$r(c_{G},c^{\prime}_{B})=(r_{GG}(c),r_{GB}(c),r_{BG}(c),r^{\prime})$. (CO8)
The good nodes know $\Lambda$, and an upper bound on the cardinalities of the
${\cal{M}}_{i}$’s, but do not know the values of the vectors $r(c)$ for any
$c\in{\cal{C}}$. (CO9) The assumption that the link-rate vector $r(c)$ does
not change with time implicitly assumes that nodes are not mobile to any
significant extent. We comment further about this assumption in Section VII,
Clock model: (CL1) Each good node $i$ has a local continuous-time clock that
it initializes to zero when it turns on. Its time $\tau^{i}(t)$ is affine with
respect to some reference time $t\geq 0$, i.e., $\tau^{i}(t)=a_{i}t+b_{i}$
where $a_{i}$ and $b_{i}$ are called the skew and offset respectively. Wlog,
the time $t$ above and in (N4) is taken equal to the clock time of the first
good node to turn on. (CL3) Denoting the relative skew and offset between
nodes $i$ and $j$ by $a_{ij}:=\frac{a_{i}}{a_{j}}$ and
$b_{ij}:=b_{i}-a_{ij}b_{j}$, node $i$’s time with respect to node $j$’s time
$s$ is $\tau^{i}_{j}(s)=a_{ij}s+b_{ij}$. We assume $0<a_{ij}\leq a_{max}$. As
a corollary of (N4,CL1,CL3), $|b_{ij}|\leq a_{max}U_{0}$, since
$\tau^{i}(U_{0})\geq 0$. (CL4) The good nodes do not know their skew or offset
a priori. (CL5) Finally, due to its digital processor, a good node $i$ can
only observe a quantized version of its continuous-time local clock
$\tau^{i}(t)$.
Cryptographic capabilities: (CR1) Each node is assigned a public key and a
private key; information encrypted by a private key can only be decrypted with
the corresponding public key. The private key is never revealed by a good node
to any other node. Possession of a public key does not enable an attacker to
forge, alter, or tamper with an encrypted packet generated with the
corresponding private key. The good nodes encrypt all their transmissions.
(CR2) Each node possesses the public key of a central authority. (CR3) Each
node possesses an identity certificate, signed by the central authority,
containing node $i$’s public key and ID number. The certificate binds node
$i$’s public key to its identity. (CR4) Each node possesses a list of all the
other $n$ node IDs.
## III The Main Results
Each time that the good nodes make a certain choice $c_{G}$, the bad nodes
could respond with some choice drawn only from a certain subset
${\cal{C}}_{B,c_{G}}\subseteq{\cal{C}}_{B}$. In this way they could ensure
that only the subset
${\cal{E}}:=\\{(c_{G},c_{B}):c_{G}\in{\cal{C}}_{G},c_{b}\in{\cal{C}}_{B,c_{G}}\\}$
is ever employed by the network. If so, we will say that ${\cal{E}}$ is
_enabled_ , while its complement
${\cal{D}}:={\cal{C}}\setminus\\{(c_{G},c_{B}):c_{G}\in{\cal{C}}_{G},c_{b}\in{\cal{C}}_{B,c_{G}}\\}$
is _disabled_ by the bad nodes. We will denote by $\Delta$ the set of all such
sets ${\cal{D}}$ that they have the capability to disable. For any set
${\cal{E}}$ of enabled CTVs, let ${{R(\cal{E})}}:=$
ConvexHull$(\\{r(c):c\in{\cal{E}}\\})$ _be the set of link rate-vectors
supported by_ ${\cal{E}}$, i.e., generated by time sharing over ${\cal{E}}$.
Let ${\cal{G}}({{\cal E}})$ be a directed graph over the nodes, where there is
an edge $ij$ if and only if $r_{ij}(c)>0$ for some $c\in{{R}}({{\cal E}})$.
We assume that the good nodes can communicate in a multi-hop fashion with each
other over bidirectional links at some minimal positive rate, regardless of
what the bad nodes do:
Connectedness Assumption (C): Let ${\cal{G}}^{*}:={\cal{G}}({{\cal
C}}\setminus{{\cal D}}^{*})$ be the graph resulting from the maximum set
${\cal{D}}^{*}\in\Delta$ that the bad nodes can disable. We will assume that
the good nodes are connected in the subgraph of ${\cal{G}}^{*}$ that consists
only of edges $ij$ for which both $ij$ as well as $ji$ are edges in
${\cal{G}}^{*}$.
Denoting by $P_{ij}$ the set of all paths from $i$ to $j$, the _multi-hop
capacity region_ of $n(n-1)$-dimensional end-to-end source-destination
throughput vectors is defined in the standard way as ${{C}}({\cal{E}}):=\\{x:$
For some vector $y\geq 0$ with $0\leq\sum_{p:\ell\in p}y_{p}\leq r_{\ell}$ for
some $r\in{{R(\cal{E})}}$, $x_{ij}=\sum_{p\in P_{ij}}y_{p}$ for all $1\leq
i,j\leq n,j\neq i\\}$.
We employ a utility function defined over the throughputs of any subset of
source-destination pairs of interest:
Utility function assumption (U): For any subset $S\subseteq\\{1,2,\ldots,n\\}$
and any throughput vector $x$, let $U(x,S)$ depend only on $x_{ij}$ for
$i,j\in S$. For every $S$, $U(x,S)$ is continuous and monotone increasing in
the components of $x$.
We now consider the game where the good nodes wish to maximize it for the
nodes perceived to be good, while the bad nodes wish to minimize it over all
their Byzantine behaviors. To obtain an upper bound on utility, suppose that
the bad nodes disable only the CTVs in ${\cal{D}}$ and reveal this choice to
the good nodes. Let ${\cal{E}}:={\cal{C}}\setminus{\cal{D}}$. If
${\cal{G}}({\cal{E}})$ has several strongly connected components, then, by the
connectedness assumption (C), the good nodes are all in the same component,
denoted by $F({\cal{E}})$, and thus know that the nodes outside $F({\cal{E}})$
are bad. They will therefore only consider the utility accrued as
$U(x,F({\cal{E}}))$, and maximize it over all $x\in{\cal{C}}({\cal{E}})$.
Hence an upper bound on achievable utility is
$\displaystyle\min_{\begin{subarray}{c}{{\cal{D}}\in\Delta}\end{subarray}}\hskip
5.69054pt\displaystyle\max_{\begin{subarray}{c}x\in{\cal{C}}({\cal{C}}\setminus{\cal{D}})\end{subarray}}U(x,F({\cal{C}}\setminus{\cal{D}}))$.
Our main result, elaborated on in Theorem VI.2, is:
###### Theorem III.1
Consider a network that satisfies (N), (CO), (CL), (CR), (C) and (U). Given an
arbitrary $\epsilon$, where $0<\epsilon<1$, the protocol described in Section
V ensures that all the good nodes obtain a common estimate of the component
that they are all members of, and achieves the utility
$\displaystyle(1-\epsilon)\displaystyle\min_{\begin{subarray}{c}{{\cal{D}}\in\Delta}\end{subarray}}\hskip
5.69054pt\displaystyle\max_{\begin{subarray}{c}x\in{\cal{C}}({\cal{C}}\setminus{\cal{D}})\end{subarray}}U(x,F({\cal{C}}\setminus{\cal{D}})).$
(1)
Some important consequences are the following. Normally, one would expect that
since the good nodes have to first declare their protocol and follow it, they
can only attain “max-min,” which is generally smaller than min-max. Since the
latter can be attained (arbitrarily closely), it shows firstly that the bad
nodes are unable to benefit from having a priori knowledge of the protocol.
Second, since all that the bad nodes can benefit from is deciding which sets
to disable, they are effectively limited to jamming and/or cooperating in each
CTV. Other more Byzantine behaviors are not any more effective.
The example below shows why a bad node may prefer to “conform” rather than jam
for some utility functions.
###### Example III.1
Consider the network of Figure 1. Nodes 1 and 2 are good and in close
proximity, while node 3 is bad and located far away. Consider the “fairness-
based ”utility function $U(x):=\min\\{x_{12},x_{32}\\}$. If node 3 jams, then
the connected component becomes $\\{1,2\\}$, and the good nodes proceed to
maximize only $x_{12}$, which node 3 can only slightly impinge because it is
so far away from node 2. However, if node 3 cooperates, then the connected
component is $\\{1,2,3\\}$, and the optimal solution for this “fair” utility
function is to make $x_{32}=x_{12}$. However, link 32 being weak, it requires
much more airtime than link 12, thus considerably reducing $x_{12}$.
Figure 1: Example III.1.
## IV The Outline of the Approach
The heart of the approach is to investigate different CTVs, exploiting the
fact that the operation of the network consists of invoking which such set to
use at any given instant. If a good node fails to receive a scheduled packet
transmitted during a CTV set, then that good node alerts the rest of the
network during a verification phase, and the offending CTV set is never used
again. After each such pruning the network then re-optimizes its utility over
the remaining CTVs. The decreasing sequence of remaining sets of CTVs
necessarily converges to an operational collection of CTVs, over which the
utility is optimized by time sharing. Since the set of disabled CTVs is
determinable by the network, as we show, it is the same as if it were revealed
to the good nodes a priori, which allows achievement of min-max. It also shows
why more complex Byzantine behaviors than jamming or cooperating are not any
more effective for the bad nodes.
There are however several problems that lie along the way to realizing this
scheme. First, all of the above presumes that all the nodes are good, and,
second, also that the nodes know the network topology and other parameters,
both of which are false. This leads to the challenge: How to determine the
network, while under attack from bad nodes when one does not know the network?
We present a complete protocol suite that proceeds through several phases to
achieve this end result.
After their birth, the nodes need to first discover who their neighbors are.
This requires a two-way handshake, which presents one problem already. Two
good nodes that are neighbors can successfully send packets to each other if
there are no primary (half-duplex) or secondary (collision) conflicts. To
achieve this we employ an Orthogonal MAC Code [18]. Next, the two nodes need
to update their clock parameters. After this, the nodes propagate their
neighborhood information so that everyone learns about the network topology.
This also poses some challenges when there are intermediary bad nodes. This is
addressed by a version of the Byzantine General’s algorithm of [1], by
capitalizing on connectedness assumption (C). Next, even though all the good
nodes converge to a common network view, that view may be internally
inconsistent, especially with respect to clocks. To resolve this we employ a
certain consistency check algorithm. Next, the nodes proceed to determine an
optimal schedule for time sharing over the set of CTVs that have performed
consistently from the very beginning, and execute it. However, a bad node that
has cooperated hitherto may not cooperate at this point. Hence the results of
this operational phase need to be verified, the dysfunctional CTV pruned, the
schedule re-optimized, and the procedure iterated.
The reader may wonder: Why do we even need a notion of “time”? First, without
it, we cannot even speak of throughput or thus of utility. Second, we use
local clocks to schedule transmissions and coordinate activity (as is quite
common, e.g., time-outs in MAC and transport protocols). On the other hand,
dependence on distributed synchronized clocks for coordinated activity opens
yet another avenue for bad nodes to sabotage the protocol – interfering with
the clock synchronization algorithm. Therefore, topics like scheduling, clock
synchronization, utility maximization, and security, are deeply interwoven.
Therefore one needs a holistic approach that addresses all these issues at
every stage of the operating lifetime, and guarantees overall security and
min-max optimality. This is the raison d’$\rm{\hat{e}}$tre for this paper.
## V The Phases of the Protocol Suite
The protocol suite consists of six phases: Neighbor Discovery, Network
Discovery, Consistency Check, Scheduling, Data Transfer, and Verification.
Proofs are deferred to Section VI.
We first note the necessity for a key ingredient. Even two good nodes that are
neighbors as in assumption (C) are only guaranteed to be able to successfully
send packets to each other provided one is transmitting, the other is
listening (since good nodes are half-duplex), and the remaining good nodes are
all silent. The Orthogonal MAC Code (OMC) of [18] ensures the simultaneity of
all these events, even though the clocks of different nodes have different
skews and offsets. For each pair of nodes $i,j$, it defines certain zero-one
valued functions of local time at each node, such that if $i$ transmits a
packet of duration $W$ to $j$ at that time, then the packet is successfully
received, and the delay involved in waiting for such an eventuality is never
more than a certain $T_{MAC}(W)$.
### V-A The Neighbor Discovery Phase
In this phase, each node $i$ will determine the identity and relative clock
parameters of nodes in its neighborhood ${\cal N}_{i}$, and include this data
in a mutually authenticated link certificate.
In the first two steps, each node $i$ attempts a handshake with a neighbor
node $j$ by broadcasting a probe packet $PRB_{ij}$ and waiting for an
acknowledgement $ACK_{ji}$. The probe packet contains an identity certificate
signed by a central authority. Given ${\cal N}_{i}:=\\{1,\ldots,n\\}$\$i$, an
initial candidate for the set of bidirectional neighbors of $i$ (as in (C)),
to indicate that node $i$ transmits $PRB_{ij}$ to each node $j\in{\cal N}_{i}$
via the OMC, and receives $PRB_{jj}$ from each node $j\in{\cal N}_{i}$, we use
$\text{TxRxMAC}(PRB_{i\rightarrow{\cal N}_{i}}$,$PRB_{{\cal N}_{i}\rightarrow
i})$. If a probe packet is not received from some node $j$, then $j$ is pruned
from ${\cal N}_{i}$.
Next, node $i$ transmits an acknowledgment $ACK_{ij}$ to node $j$ containing a
signed confirmation of the received probe packet $PRB_{j}$. Node $i$ also
listens for an acknowledgment $ACK_{ji}$ from node $j$. Node $i$ further
removes from ${\cal N}_{i}$ any nodes that failed to return acknowledgements.
Then node $i$ transmits to each node $j\in{\cal N}_{i}$ a pair of timing
packets $TIM^{(1)}_{i,j}$ and $TIM^{(2)}_{j,i}$ that contain the send-times
$s^{(1)}_{ij}$ and $s^{(2)}_{ij}$ respectively as recorded by its local clock
$\tau^{j}(t)$. Node $i$ also receives a corresponding pair of timing packets
$TIM^{(1)}_{i,j}$ and $TIM^{(2)}_{j,i}$ from node $j$, and records the
corresponding receive-times $r^{(1)}_{ji}$ and $r^{(2)}_{ji}$ respectively, as
measured by the local clock $\tau^{i}(t)$. Any node that fails to deliver
timing packets to node $i$ is further removed from ${\cal N}_{i}$. The timing
packets are used to estimate the relative skew $a_{ji}$ by
$\hat{a}_{ji}:=\frac{r^{(2)}_{ji}-r^{(1)}_{ji}}{s^{(2)}_{ji}-s^{(1)}_{ji}}$.
The relative skew is used at the end of the Network Discovery Phase, to
estimate a reference clock with respect to the local continuous-time clock. In
the last two steps, node $i$ creates a link certificate $LNK^{(1)}_{ij}$
containing the computed relative clock skew with respect to node $j$, and
transmits this link to node $j$ using the OMC. Node $i$ also listens for a
similar link certificate $LNK^{(2)}_{ji}$ from node $j$.
Finally, node $i$ verifies and signs the received link certificate, and
transmits the authenticated version $LNK^{(2)}_{ij}$ back to node $j$. Node
$i$ listens for a similar authenticated link certificate $LNK^{(2)}_{ji}$ from
$j$. Any nodes that fail to return link certificates are removed from the set
${\cal N}_{i}$. This set now represents the nodes in the neighborhood of node
$i$ with whom node $i$ has established mutually authenticated link
certificates. The Neighbor Discovery Phase’s pseudocode is shown in Algorithm
1.
One problem is that the algorithm must be completed in a partially coordinated
manner even though the nodes are asynchronous; the completion of any stage in
the Exponential Information Gathering (EIG) algorithm (see below) depends on
the successful completion of the previous stages by all other good nodes.
Consequently, we assign increasingly larger intervals
$S_{k}:=[t_{k},t_{k+1}),k=1,\ldots 6,$ to each successive protocol stage; see
Section VI.
Algorithm 1 The Neighbor Discovery Phase
procedure NeighborDiscovery
${\cal N}_{i}:=\\{1,\ldots,n\\}\setminus i$
while $t\in S_{1}$ do
TxRxMAC($PRB_{i\rightarrow{\cal N}_{i}}$,$PRB_{{\cal N}_{i}\rightarrow i}$)
Update(${\cal N}_{i}$)
end while
while $t\in S_{2}$ do
TxRxMAC($ACK_{i\rightarrow{\cal N}_{i}}$,$ACK_{{\cal N}_{i}\rightarrow i}$)
end while
while $t\in S_{3}$ do
TxRxMAC($TIM^{(1)}_{i\rightarrow{\cal N}_{i}}$,$TIM^{(1)}_{{\cal
N}_{i}\rightarrow i}$)
Update(${\cal N}_{i}$)
end while
while $t\in S_{4}$ do
TxRxMAC($TIM^{(2)}_{i\rightarrow{\cal N}_{i}}$,$TIM^{(2)}_{{\cal
N}_{i}\rightarrow i}$)
Update(${\cal N}_{i}$)
end while
while $t\in S_{5}$ do
TxRxMAC($LNK^{(1)}_{i\rightarrow{\cal N}_{i}}$,$LNK^{(1)}_{{\cal
N}_{i}\rightarrow i}$)
Update(${\cal N}_{i}$)
end while
while $t\in S_{6}$ do
TxRxMAC($LNK^{(2)}_{i\rightarrow{\cal N}_{i}}$,$LNK^{(2)}_{{\cal
N}_{i}\rightarrow i}$)
Update(${\cal N}_{i}$)
end while
end procedure
### V-B The Network Discovery Phase
The purpose of this Phase is to allow the good nodes to obtain a _common_ view
of the network topology and _consistent_ estimates of all clock parameters. To
accomplish this, the good nodes must disseminate their lists of neighbors to
all nodes, so that all can decide on the same topology view. However good
nodes do not know a priori which nodes are bad, and so bad nodes can
selectively drop lists or introduce false lists to prevent consensus. We
resolve this by using a version of the Byzantine General’s algorithm of [1],
requiring an EIG tree data structure. Let $T_{i}$ denote node $i$’s EIG tree,
which by construction has depth $n$. The root of $T_{i}$ is labelled with node
$i$’s neighborhood, i.e., the nodes in ${\cal N}_{i}$ and the corresponding
collection of link certificates. First node $i$ transmits to every node
$j\in{\cal N}_{j}$ in its neighborhood, the list of nodes in ${\cal N}_{i}$
and corresponding link certificates, while receiving similar lists from each
node in ${\cal N}_{j}$. Node $i$ updates its EIG tree with the newly received
lists from its neighbors, by assigning each received list to a unique child
vertex of the root of $T_{i}$. Node $i$ then transmits the set of level 1
vertices of $T_{i}$ to every node in its neighborhood, receiving a set of
level 1 vertcies from each neighbor in turn. The EIG tree $T_{i}$ is updated
again. This process continues through all $n$ levels of the EIG tree.
The notation $T^{(k)}_{i}$ in Algorithm 2 indicates the $k$-level vertices of
the EIG tree $T_{i}$. The notation $\text{TxRxMAC}(T^{(k)}_{i\rightarrow{\cal
N}_{i}},T^{(k)}_{{\cal N}_{i}\rightarrow i})$ indicates that, using the OMC,
node $i$ transmits $T^{(k)}_{i}$ to each node $j\in{\cal N}_{i}$, and receives
$T^{(k)}_{j}$ from each node $j\in{\cal N}_{i}$.
We use $\text{UPDATE}(T_{i})$, to update the EIG tree $T_{i}$ after the
arrival of new information, and the procedure $\text{DECIDE}(T_{i})$ to infer
the network topology based on the EIG tree. The $n$-stage EIG algorithm
guarantees that if the subgraph of good nodes is connected, then each good
node will decide on the same topological view.
Algorithm 2 The EIG Byzantine General’s Algorithm
procedure EIGByzMAC(${\cal N}_{i}$)
$T^{(0)}_{i}:={\cal N}_{i}$
for $k=1,\ldots n$ do
while $t\in S_{6+k}$ do
TxRxMAC($T^{(k)}_{i\rightarrow{\cal N}_{i}},T^{(k)}_{{\cal N}_{i}\rightarrow
i}$)
Update($T_{i}$)
end while
end for
Decide($T_{i}$)
end procedure
### V-C The Consistency Check Phase
Unfortunately, a fundamental difficulty is that malicious nodes along a path
$1,\ldots,n$ may have generated false time stamps in the Neighbor Discovery
Phase, and thus corrupted the measured relative skews between adjacent nodes.
There may be several connecting paths infiltrated by bad nodes that thereby
generate different values for the relative skew. It is impossible to determine
the correct path from the relative skews alone. Every pair of such
inconsistent paths corresponds to an inconsistent cycle in which the skew
product is not equal to one. We use an algorithm called Consistency Check to
identify the path that generated the correct relative skew.
Consistency Check works by circling a timing packet around every cycle in
which the skew product differs from one by more than $\epsilon_{a}$, a desired
maximum skew error. At the conclusion of the test, at least one link with a
malicious endpoint will be removed from the cycle, eliminating a connecting
path. During the test, each node in such a cycle is obliged to append a
receive time-stamp and a send time-stamp generated by the local clock before
forwarding the packet to the next node. These time-stamps must satisfy a delay
bound condition; the send time and receive time cannot differ by more than 1
clock count. A node fails the consistency check otherwise, or if its time
stamps do not agree with its declared relative skew. The key idea is that if
the test starts after a sufficiently large amount of time has elapsed, the
clock estimates based on faulty relative skews will have diverged so
extensively from the actual clocks that at least one malicious node in the
cycle will find it impossible to generate time-stamps that are consistent with
its declared relative clock skew and satisfy the delay bound condition (all
proofs are in Section VI):.
###### Theorem V.1
Let $T_{j}$ be the start-time of the Consistency Check for the $jth$
inconsistent cycle, consisting of nodes $i_{1},\ldots,i_{m}$. At least one
malicious node in cycle $j$ will violate a consistency check condition, if
$T_{j}>\frac{\hat{a}_{i_{m},i^{*}}(m+1)K+\epsilon_{b}}{\epsilon_{a}}$ where
$i^{*}$ is the node with the smallest skew product $\hat{a}_{i^{*},i_{1}}$.
Algorithm 3 depicts Consistency Check. Given a cycle, $j$, $k$ and $m$ denote
nodes that follow and precede node $i$ respectively in the cycle. If node $i$
is the leader of the cycle, i..e., the node with smallest ID, then node $i$
initiates the timing packet that traverses the cycle and transmits it to node
$k$. Otherwise, node $i$ waits for the timing packet to arrive from node $m$
before forwarding it to node $k$.
Algorithm 3 Consistency Check Algorithm at Node $i$
procedure ConsistencyCheck
$START:=\frac{(n+1)(a_{max})^{n+1}+(n+1)(a_{max})^{n+1}U_{0}}{\epsilon_{a}}$
for each cycle $C_{j}$ do
$k=$Next($C_{j}$)
$m=$Prev($C_{j}$)
if $i$=Leader($C_{j}$) and $t\geq START$ then
Transmit($TIM_{i\rightarrow k}$)
else if $i\in C_{j}$ then
Receive($TIM_{m\rightarrow i}$)
Transmit($TIM_{i\rightarrow k}$)
end if
end for
end procedure
After all inconsistent cycles have been tested, each node $i$ disseminates the
set of all timing packets ${\cal T}_{i}$ it received to other nodes. The EIG
algorithm is used to ensure a common view of the timing packets generated.
Each node removes from the topology any link whose endpoints generate time-
stamps inconsistent with its declared relative skew or violated the delay
bound. The complete Phase is shown in Algorithm 4.
Algorithm 4 The Network Discovery Phase at Node $i$
procedure NetworkDiscovery
EIGByzMAC(${\cal N}_{i}$)
ConsistencyCheck
EIGByzMAC(${\cal T}_{i}$)
end procedure
At the conclusion of Network Discovery Phase node $i$ shares a common view of
the network topology with all other good nodes. As a result, the network can
designate the node with smallest ID as the _reference clock_. Furthermore,
each node $i$ has an estimate of the reference clock $\tau^{r}_{i}(t)$ with
respect to its local clock $t$ using the formula
$\hat{\tau}^{r}_{i}(t):=\hat{a}_{ri}t$, where estimated $\hat{a}_{ri}$ and
actual relative skews $a_{ri}$ differ by at most $\epsilon_{a}$.
### V-D The Scheduling Phase
In the Scheduling Phase the good nodes in the network obtain a common schedule
governing the transmission and reception of data packets. A “schedule” is
simply a sequence of CTVs, each with specified start and end times. Each node
$i$ divides the Data Transfer Phase into time-slots, and assigns a CTV to each
time-slot so that the resulting throughput vector is utility optimal. All the
good nodes independently arrive at the same schedule since they independently
optimize the same utility function over the same ${\cal C}$ (ties broken
lexicographically).
Since the good nodes must conform to a common schedule, each node $i$
generates a local estimate of the reference clock $\hat{\tau}^{r}_{i}(t)$ with
respect to its local clock $t$, as described in the Network Discovery Phase.
However, this estimate may not be perfectly accurate; some of the nodes on a
path along which relative skew is estimated may be malicious and can introduce
an error of at most $\epsilon_{a}$ into the computed relative skew. To address
this, the time-slots are separated by a dead-time of size $D$, where given any
pair of nodes $(i,j)$, $D$ is chosen to satisfy
$|\hat{\tau}^{r}_{i}(t)-\hat{\tau}^{r}_{j}(\tau^{j}_{i}(t))|\leq D$.
Finally, $n^{2}(n-1)$ time-slots are enough to guarantee that every pair of
nodes can communicate once in either direction, via multihop routing, during
Data Transfer Phase. The algorithm $\text{UtilityMaximization}({\cal C})$ for
the Scheduling Phase is depicted in Algorithm 5. At the end of Scheduling
Phase, node $i$ shares a common utility maximizing schedule with other good
nodes.
Algorithm 5 The Scheduling Phase at Node $i$
procedure Scheduling
UtilityMaximization(${\cal C}$)
end procedure
### V-E The Data Transfer Phase
In this Phase the nodes exchange data packets using the generated schedule. It
is divided into time-slots, with each assigned a CTV, a rate vector, and set
of packets for each transmitter in the set. To prevent collisions resulting
from two nodes assigning themselves to different time slots due to timing
error, node $i$ begins transmission $D$ time-units after the start of the
time-slot. The transmitted packet is then guaranteed to arrive at the receiver
in the same time slot, for appropriate choice of $D$ and time-slot size
$B_{slot}$.
Algorithm 7 defines this phase, with $m_{k}$ denoting a message to be
transmitted or received by node $i$ in the $k_{th}$ slot, $T_{start}$ the
start time of the phase measured by the local estimate of the reference clock
$\hat{\tau}^{r}_{i}(t))$, $S_{k}=[t_{k},t_{k+1}),k=1,\ldots,N$ the time-slots
of the phase with $N=n^{2}(n-1)$, $t_{1}=T_{start}$, and
$t_{k+1}:=t_{k}+B_{slot}+2D$, and $TX(k)$ and $RX(k)$ the CTV, and receiving
nodes during slot $k$.
Algorithm 6 The Data Transfer Phase at Node $i$
procedure DataTransfer($T_{start}$)
for k=1,…,N do
if $t\in S_{k}$ and $t\geq t_{k}+D$ and $i\in TX(k)$ then
Transmit($m_{k}$)
else if $t\in S_{k}$ and $i\in RX(k)$ then
Receive($m_{k}$)
end if
end for
end procedure
### V-F The Verification Phase
However, malicious nodes may not cooperate in the Data Transfer Phase. So
whenever a scheduled packet fails to arrive at node $j$, it adds the offending
CTV and associated packet number to a list, and disseminates the list in the
Verification Phase using the EIG Byzantine General’s algorithm. These CTVs are
then permanently further pruned from the collection of feasible CTVs. With
$L_{k}$ denoting the list that failed during the $k$th iteration of the Data
Transfer Phase, the set ${\cal C}_{k}$ of feasible CTVs during the $k$th
iteration of the Scheduling Phase is updated to ${\cal C}_{k+1}={\cal
C}_{k}\setminus L_{k}$ in Algorithm 7.
All communication can be scheduled into slots separated by a dead-time of
$2D$. Within each of the $n$ stages of the EIG Byzantine General’s algorithm,
there are $n(n-1)$ pairs of nodes that may communicate, and at most $n$ nodes
on the connecting path. Therefore, the total number of time slots required is
$n^{3}(n-1)$.
At the conclusion of the phase, the good nodes again share a common view of
the set of feasible CTVs for the next iteration of the Scheduling Phase.
Algorithm 7 The Verification Phase at Node $i$
procedure Verification
EIGByz($L_{k}$)
Update(${\cal C}_{k+1}$)
end procedure
### V-G The Steady State
The network cycles through Scheduling, Data Transfer, and Verification Phases
for $n_{iter}$ iterations. Eventually, by finiteness, it converges to a set of
CTVs, and a utility-maximizing schedule over it. The overall protocol is in
Algorithm 8.
Algorithm 8 The Complete Protocol
NeighborDiscovery
NetworkDiscovery
for $k=1,\ldots,n_{iter}$ do
Scheduling(${\cal C}_{k}$)
DataTransfer(t)
Verification
end for
## VI Feasibility of Protocol and Optimality Proof
For the distributed wireless nodes to exchange data over the network, they
must not only have the same topological view, in order to independently arrive
at a common schedule, but they must also have a consistent view of a reference
clock so that any activity will conform to this common schedule. For this, we
consider the consistency check algorithm of Section V-C.
Consider a chain network $1,\ldots,n$, where the endpoints, nodes $1$ and $n$
are good, and the intermediate nodes $2,\ldots,n-1$ are bad. Note that this
network can also be reduced to a cycle of size $n-1$ by making both endpoints
the same node. We assume that the two good endpoints do not know if any of the
intermediate nodes are bad.
Now suppose that each pair of adjacent nodes $(i,i-1)$ for $i=2,\ldots,n$ has
declared a set of relative skews and offsets
$\\{\hat{a}_{i,i-1},\hat{b}_{i,i-1}\\}$, and that each node in the chain knows
this set. The two good nodes wish to determine whether the declared skews are
accurate, i.e., whether $a_{n,1}=\prod^{n}_{i=2}\hat{a}_{i,i-1}$.
Unfortunately, the good nodes have no way of directly measuring $a_{n,1}$. The
estimate of $a_{n,1}$ is obtained from the skew product itself, which is the
very quantity that needs to be verified. So, instead, the good nodes carry out
the consistency check described earlier. After waiting a sufficiently long
time, node $1$ initiates a timing packet that traverses the chain from left to
right. Each node in the chain is obligated to forward the packet after
appending receive and time-stamps that satisfy the skew consistency and delay
bound conditions.
In order to defeat this test, the bad nodes, having collectively declared a
false set of relative skews and offsets, must support two sets of clocks for
each node $i\in\\{2,\ldots,n\\}$: a “left” clock $\tau^{i,l}(t)$ to generate
receive time-stamps, and a “right” clock $\tau^{i,r}(t)$ to generate send-time
stamps. Unlike the clocks of the good nodes, the left and right clocks of the
bad nodes need not be affine with respect to the global reference clock. In
fact, the bad nodes are free to jointly select any set of clocks
$\\{\tau^{i,l}(t),\tau^{i,r}(t),\forall i=2,\ldots,n-1\\}$ that are arbitrary
functions of $t$, a much larger set than the affine clocks being emulated.
However, we will show that if node $1$ waits sufficiently long enough, there
is no set of clocks $\\{\tau^{i,l}(t),\tau^{i,r}(t),i=2,\ldots,n-1\\}$ that
can generate time-stamps which satisfy both conditions of the consistency
check.
Let $r_{i,i-1}$ and $s_{i,i+1}$ denote the receive and send time-stamps
generated by a bad node $i$ with respect to the left and right clocks
$\tau^{i,l}(t)$ and $\tau^{i,r}(t)$ respectively. Let $t_{i,l}$ and $t_{i,r}$
denote the time with respect to the global reference clock at which the
receive and send time-stamps are generated at node $i$. We have
$r_{i-1,i}:=\tau^{i,l}(t_{i,l})$ and $s_{i,i+1}:=\tau^{i,r}(t_{i,r})$. Let
$t_{1}$ and $t_{n}$ denote the time with respect to the global reference clock
at which the timing packet was transmitted by node $1$ and received by node
$n$ respectively. We have
$s_{1,2}:=\tau^{1}(t_{1}),_{n-1,n}:=\tau^{n}(t_{n})$. To simplify notation we
will define left and right clocks at the endpoints so that
$t_{1,r}:=t_{1},t_{n,l}:=t_{n}$ and
$\tau^{1,r}(t_{1,r}):=\tau^{1}(t_{1}),\tau^{n,l}(t_{n,l}):=\tau^{n}(t_{n})$.
In order to prove that both conditions of the consistency check cannot be
satisifed by any set of clocks
$\\{\tau^{i,l}(t),\tau^{i,r}(t),i=2,\ldots,n-1\\}$, we will assume that the
first condition is satisfied, and show that second must fail. Therefore, the
clocks must satisfy:
$\displaystyle\tau^{i,l}(t_{i,l})$
$\displaystyle=a_{i,i-1}\tau^{i-1,r}(t_{i-1,r})+b_{i,i-1}\mbox{ for }i\leq
2\leq n.$ (2)
In addition, by virtue of causality, we also have:
$\displaystyle\tau^{i,l}(t_{i,l})$ $\displaystyle\leq\tau^{i,r}(t_{i,r}).$ (3)
We prove that delay bound condition must be violated if node $1$ waits for a
sufficiently large period of time before before initiating the timing packet,
i.e., if $\tau^{1}(t_{1})$ is sufficiently large, then for some $i$, we have
$\tau^{i,r}(t_{i,r})-\tau^{i,l}(t_{i,l})>K$. More precisely, we show
$\sum^{n-1}_{i=2}\left(\tau^{i,r}(t_{i,r})-\tau^{i,l}(t_{i,l})\right)>nK$,
which implies that some node has violated delay bound condition.
The sum $\sum^{n-1}_{i=2}\left(\tau^{i,r}(t_{i,r})-\tau^{i,l}(t_{i,l})\right)$
cannot be directly evaluated because the left and right clocks
$\\{\tau^{i,l}(t),\tau^{i,r}(t)\\}$ are arbitrary functions of $t$. However,
we have the following equality by repeated addition and subtraction
$\tau^{n,l}(t_{n,l})=\tau^{1,r}(t_{1,r})+\sum^{n}_{i=2}\left(\tau^{i,l}(t_{i,l})-\tau^{i-1,r}(t_{i-1,r})\right)=\sum^{n-1}_{i=2}\left(\tau^{i,l}(t_{i,l})-\tau^{i-1,r}(t_{i-1,r})\right)=\tau^{1,r}(t_{1,r})+S_{1}+S_{2}$,
where
$S_{1}:=\sum^{n}_{i=2}\left(\tau^{i,l}(t_{i,l})-\tau^{i-1,r}(t_{i-1,r})\right),S_{2}:=\sum^{n-1}_{i=2}\left(\tau^{i,l}(t_{i,l})-\tau^{i-1,r}(t_{i-1,r})\right)$.
The value $S_{2}$ is the sum of the forwarding delays. We will use (2) and (3)
to obtain an upper bound on $S_{1}$. Inserting this upper bound and using the
fact that $\tau^{n,l}(t)$ and $\tau^{1,r}(t)$ are both affine functions of
$t$, will allow us to obtain a lower bound on $S_{2}$. The proof will then
follow easily. We now obtain an upper bound on $S_{1}$ when the forward skew
product $\prod^{j}_{i=2}\hat{a}_{i,i-1}\geq 1$ for all $j\geq 2$.
###### Lemma VI.1
Suppose $\prod^{j}_{i=2}a_{i,i-1}\geq 1$ for $2\leq i\leq n$. Then
$\sum^{n}_{i=2}(\tau^{i,l}(t_{i,l})-\tau^{i-1,r}(t_{i-1,r}))\leq\left(\frac{\hat{a}_{n,1}-1}{\hat{a}_{n,1}}\right)\tau^{n,l}(t_{n,l})\sum^{n}_{i=2}\frac{\hat{b}_{i,i-1}}{\hat{a}_{i,1}}$.
###### Proof:
We have by definition
$\tau^{n+1,l}(t_{n,l}):=\hat{a}_{n+1,n}\tau^{n,r}(t_{n,r})+\hat{b}_{n+1,n}$.
For $n=2$, we have
$\tau^{2,l}(t_{2,l})-\tau^{1,r}(t_{1,r})=\left(\frac{a_{2,1}-1}{a_{2,1}}\right)\tau^{2,l}(t_{2,l})+\frac{b_{2,1}}{a_{2,1}}$.
Now assume the lemma is true for $n$. We will show that it also holds for
$n+1$:
$\sum^{n+1}_{i=2}(\tau^{i,l}(t_{i,l})-\tau^{i-1,r}(t_{i-1,r}))=\sum^{n}_{i=2}(\tau^{i,l}(t_{i,l})-\tau^{i-1,r}(t_{i-1,r}))+\tau^{n+1,l}(t_{n+1,l})-\tau^{n,r}(t_{n,r})\leq\left(\frac{\hat{a}_{n,1}-1}{\hat{a}_{n,1}}\right)\tau^{n,l}(t_{n,l})+\sum^{n}_{i=2}\frac{\hat{b}_{i,i-1}}{\hat{a}_{i,1}}+\tau^{n+1,l}(t_{n+1,l})-\tau^{n,r}(t_{n,r})\leq\left(\frac{\hat{a}_{n,1}-1}{\hat{a}_{n,1}}\right)\tau^{n,r}(t_{n,r})+\sum^{n}_{i=2}\frac{\hat{b}_{i,i-1}}{\hat{a}_{i,1}}+\tau^{n+1,l}(t_{n+1,l})-\tau^{n,r}(t_{n,r})=\left(\frac{\hat{a}_{n,1}-1}{\hat{a}_{n,1}}\right)\left(\frac{\tau^{n+1,l}(t_{n+1,l})-\hat{b}_{n+1,n}}{\hat{a}_{n+1,n}}\right)+\sum^{n}_{i=2}\frac{\hat{b}_{i,i-1}}{\hat{a}_{i,1}}+\left(\frac{\hat{a}_{n+1,n}-1}{\hat{a}_{n+1,n}}\right)\hat{\tau}^{n+1,l}(t_{n+1,l})+\frac{\hat{b}_{n+1,n}}{\hat{a}_{n+1,n}}=\left(\frac{\hat{a}_{n+1,1}-1}{\hat{a}_{n+1,1}}\right)\tau^{n+1,l}(t_{n+1,l})+\sum^{n+1}_{i=2}\frac{\hat{b}_{i,i-1}}{\hat{a}_{i,1}}$,
which follow from the induction hypothesis above in the Lemma statement, and
the fact that $\tau^{n,r}(t_{n,r})\geq\tau^{n,l}(t_{n,l})$ and $a_{i,1}\geq 1$
for all $2\leq i\leq n+1$ (that is, the coefficient $\hat{a}_{i,1}-1$ is
negative). ∎
We next bound $S_{1}$ in the special case when the reverse skew product
$\prod^{j}_{i=1}\hat{a}_{n-(i-1),n-i}\leq 1$ for all $j\geq 1$.
###### Lemma VI.2
Suppose $\prod^{j}_{i=1}a_{n-(i-1),n-i}\leq 1$ for $2\leq j\leq n-1$. Then
$\sum^{j}_{i=1}(\tau^{n-(i-1),l}(t_{n-(i-1),l})-\tau^{n-i,r}(t_{n-i,r}))$
$\leq\left(\hat{a}_{n,n-j}-1\right)\tau^{n-j,r}(t_{n-j,r})+\hat{b}_{n,n-1}+\sum^{n-1}_{i=n-j+1}\hat{a}_{n,i}\hat{b}_{i,i-1}$.
###### Proof:
We have by definition
$\tau^{n-(k-1),l}(t_{n-(k-1),l}):=\hat{a}_{n-(k-1),n-k}\tau^{n-k,r}(t_{n-k,r})+\hat{b}_{n-(k-1),n-k}$.
For $j=1$,
$\tau^{n,l}(t_{n,l})-\tau^{n-1,r}(t_{n-1,r})=(a_{n,n-1}-1)\tau^{n-1,r}(t_{n-1,r})$.
Now assume the Lemma holds for $j$. We will show that it must hold for $j+1$:
$\sum^{j+1}_{k=1}(\tau^{n-(k-1),l}(t_{n-(k-1),l})-\tau^{n-k,r}(t_{n-k,r}))=\sum^{j}_{k=1}(\tau^{n-(k-1),l}(t_{n-(k-1),l})-\tau^{n-k,r}(t_{n-k,r}))+\tau^{n-j,l}(t_{n-j,l})-\tau^{n-(j+1),r}(t_{n-(j+1),r})\leq(\hat{a}_{n,n-j}-1)\tau^{n-j,r}(t_{n-j,r})+\hat{b}_{n,n-1}+\sum^{n-1}_{k=n-j+1}\hat{a}_{n,k}\hat{b}_{k,k-1}+\tau^{n-j,l}(t_{n-j,l})-\tau^{n-(j+1),r}(t_{n-(j+1),r})\leq(\hat{a}_{n,n-j}-1)\tau^{n-j,l}(t_{n-j,l})+\hat{b}_{n,n-1}+\sum^{n-1}_{k=n-j+1}\hat{a}_{n,k}\hat{b}_{k,k-1}+\tau^{n-j,l}(t_{n-j,l})-\tau^{n-(j+1),r}(t_{n-(j+1),r})\leq(\hat{a}_{n,n-(j+1)}-1)\tau^{n-(j+1),r}(t_{n-(j+1),r})+\hat{b}_{n,n-1}+\sum^{n-1}_{k=n-j}\hat{a}_{n,k}\hat{b}_{k,k-1}$.
The above follow from induction hypothesis in Lemma VI.2, since
$\tau^{i,l}(t_{i,l})\leq\tau^{i,r}(t_{i,r})$ and $\hat{a}_{n,n-j}\leq 1$ for
$1\leq j\leq n-1$ (that is, the coefficient $\hat{a}_{n,n-j}-1$ is negative),
and from substitution into $\tau^{n-j,l}(t_{n-j,l})$ and simplification. ∎
We will combine both special cases in Lemma VI.1 and Lemma VI.2 to obtain an
upper bound on $S_{1}$. First we define $i^{*}$ as the node with the smallest
skew product $\hat{a}_{i^{*},1}$ in the chain network, that is less than one.
That is, $\hat{a}_{i^{*},1}=\displaystyle\min_{k}\hat{a}_{k,1}$ and
$\hat{a}_{i^{*},1}\leq 1$. If no such node exists, set $i^{*}=1$.
Now we consider an arbitrary set of skews
$\\{\hat{a}_{i,i-1},i=2,\ldots,n\\}$. Next we show that if $i^{*}\geq 2$ then
the forward skew product starting from $i^{*}$ is greater than 1, and the
reverse skew product starting from $i^{*}-1$ is always less than one.
###### Lemma VI.3
If $i^{*}\geq 2$ then $\hat{a}_{j,i^{*}}\geq 1$ for $i^{*}+1\leq j\leq n$ and
$\hat{a}_{i^{*},i^{*}-k+1}\leq 1$ for $1\leq k\leq i^{*}$. Otherwise,
$\hat{a}_{j,1}\geq 1$ for $2\leq j\leq n$.
###### Proof:
Consider $i^{*}\geq 2$, and suppose the first part of the assertion is false.
I.e., for some $j^{\prime}$, $\hat{a}_{j^{\prime}i^{*}}<1$. It follows that
$\hat{a}_{j^{\prime}1}=\hat{a}_{j^{\prime}i^{*}}\hat{a}_{i^{*}1}\leq\hat{a}_{i^{*}1}$.
But then $j^{\prime}$ is a node with a smaller skew product $\hat{a}_{j1}$
than node $i^{*}$, which contradicts the definition of $i^{*}$. Now suppose
that the second part of the assertion is false. I.e., for some $j^{\prime}$ we
have $\hat{a}_{i^{*}j^{\prime}}>1$. It follows that
$\hat{a}_{i^{*}1}=\hat{a}_{i^{*}j^{\prime}}\hat{a}_{j^{\prime}1}\geq\hat{a}_{j^{\prime}1}$.
But then $j^{\prime}$ is a node with a smaller skew product than node $i^{*}$,
which again contradicts the definition of $i^{*}$. Now consider the case when
$i^{*}=1$. Then by definition of $i^{*}$ it follows that $\hat{a}_{j1}\geq 1$
for all $2\leq j\leq n$. ∎
We now obtain an upper bound on $S_{1}$ for arbitrary skews.
###### Lemma VI.4
Suppose $i^{*}\geq 2$. We have the following inequality:
$\sum^{n}_{j=2}\tau^{j,l}(t_{j,l})-\tau^{j-1,r}(t_{j-1,r})\leq(\hat{a}_{i^{*},1}-1)\tau^{1,r}(t_{1,r})$
$+\left(\frac{\hat{a}_{n,i^{*}}-1}{\hat{a}_{n,i^{*}}}\right)\tau^{n,l}(t_{n,l})+\frac{\hat{b}_{n,1}}{\hat{a}_{n,i^{*}}}$.
###### Proof:
$\sum^{n}_{j=2}\tau^{j,l}(t_{j,l})-\tau^{j-1,r}(t_{j-1,r})$
$\sum^{i^{*}}_{j=2}\tau^{j,l}(t_{j,l})-\tau^{j-1,r}(t_{j-1,r})+\sum^{n}_{j=i^{*}+1}\tau^{j,l}(t_{j,l})-\tau^{j-1,r}(t_{j-1,r})$
$=(\hat{a}_{i^{*},1}-1)\tau^{1,r}(t_{1,r})+\hat{b}_{i^{*},i^{*}-1}$
$+\sum^{i^{*}}_{j=2}\hat{a}_{i^{*},j}\hat{b}_{j,j-1}+\left(\frac{\hat{a}_{n,i^{*}}-1}{\hat{a}_{n,i^{*}}}\right)\tau^{n,l}(t_{n,l})+\sum^{n}_{i=i^{*}+1}\frac{\hat{b}_{i,i-1}}{\hat{a}_{i,i^{*}}}$
$=(\hat{a}_{i^{*},1}-1)\tau^{1,r}(t_{1,r})+\left(\frac{\hat{a}_{n,i^{*}}-1}{\hat{a}_{n,i^{*}}}\right)\tau^{n,l}(t_{n,l})$
$+\sum^{n}_{j=2}\frac{\hat{a}_{n,j}\hat{b}_{j,j-1}}{\hat{a}_{n,i^{*}}}$
$=(\hat{a}_{i^{*},1}-1)\tau^{1,r}(t_{1,r})+\left(\frac{\hat{a}_{n,i^{*}}-1}{\hat{a}_{n,i^{*}}}\right)\tau^{n,l}(t_{n,l})+\frac{\hat{b}_{n,1}}{\hat{a}_{n,i^{*}}}.$which
follow by applying Lemma VI.2 and Lemma VI.1, by multiplying the terms in each
summation by $\frac{\hat{a}_{n,i^{*}}}{\hat{a}_{n,i^{*}}}$ and simplifying,
and from the definitions of $\hat{b}_{ij}$ and $\hat{d}^{(i)}_{ji}$. ∎
Now that we have an upper bound on $S_{1}$, we can obtain a lower bound on
$S_{2}$, the sum of the forwarding delays.
###### Lemma VI.5
The sum of forwarding delays in the chain network satisfies:
$\sum^{n-1}_{j=2}\left(\tau^{j,l}(t_{j,l})-\tau^{j,r}(t_{j,r})\right)\geq\frac{(a_{n,1}-\hat{a}_{n,1})}{\hat{a}_{n,i^{*}}}\tau^{1,r}(t_{1,r})+\frac{(b_{n,1}-\hat{b}_{n,1})}{\hat{a}_{n,i^{*}}}$.
###### Proof:
$\sum^{n-1}_{j=2}\left(\tau^{j,l}(t_{j,l})-\tau^{j,r}(t_{j,r})\right)$
$=\tau^{n,l}(t_{n,l})-\tau^{n,r}(t_{n,r})-\sum^{n}_{j=2}\tau^{j,l}(t_{j,l})-\tau^{j-1,r}(t_{j-1,r})$
$\geq\tau^{n,l}(t_{n,l})-\tau^{n,r}(t_{n,r})-(\hat{a}_{i^{*},1}-1)\tau^{1,r}(t_{1,r})$
$-\left(\frac{\hat{a}_{n,i^{*}}-1}{\hat{a}_{n,i^{*}}}\right)\tau^{n,l}(t_{n,l})-\frac{\hat{b}_{n,1}}{\hat{a}_{n,i^{*}}}$
$=\frac{\tau^{n,l}(t_{n,l})}{\hat{a}_{n,i^{*}}}-\hat{a}_{i^{*},1}\tau^{1,r}(t_{1,r})-\frac{\hat{b}_{n,1}}{\hat{a}_{n,i^{*}}}$
$\geq\frac{\tau^{n,l}(t_{1,r})}{\hat{a}_{n,i^{*}}}-\hat{a}_{i^{*},1}\tau^{1,r}(t_{1,r})-\frac{\hat{b}_{n,1}}{\hat{a}_{n,i^{*}}}$
$=\frac{a_{n,1}\tau^{1,r}(t_{1,r})+b_{n,1}}{\hat{a}_{n,i^{*}}}-\hat{a}_{i^{*},1}\tau^{1,r}(t_{1,r})-\frac{\hat{b}_{n,1}}{\hat{a}_{n,i^{*}}}$
$=\frac{(a_{n,1}-\hat{a}_{n,1})}{\hat{a}_{n,i^{*}}}\tau^{1,r}(t_{1,r})+\frac{(b_{n,1}-\hat{b}_{n,1})}{\hat{a}_{n,i^{*}}}$,
which follow by noting from repeated addition and subtraction that
$\tau^{n,l}(t_{n,l})=\tau^{1,r}(t_{1,r})+\sum^{n}_{j=2}\left(\tau^{j,l}(t_{j,l})-\tau^{j-1,r}(t_{j-1,r})\right)+\sum^{n-1}_{j=2}\left(\tau^{j,l}(t_{j,l})-\tau^{j,r}(t_{j,r})\right)$,
by applying Lemma VI.4, because $t_{n,l}\geq t_{1,r}$ since node $n$ could not
have received the timing packet before node $1$ transmitted it, and since node
$n$’s clock is relatively affine with respect to node $1$’s clock. ∎
We now complete the proof of consistency check for a chain network. We show
that if the start time of the consistency check is sufficiently large, and the
left and right clocks $\\{\tau^{i,l}(t_{i,l}),\tau^{i,r}(t_{i,r})\\}$ satisfy
the parameter consistency condition, then at least one node will violate delay
bound condition. Hence there are no left and right clocks that can pass both
conditions of consistency check if start time is large.
###### Proof:
We assume node $1$ is a good node. Now
$\frac{(a_{n,1}-\hat{a}_{n,1})}{\hat{a}_{n,i^{*}}}\tau^{1,r}(t_{1,r})+\frac{(b_{n,1}-\hat{b}_{n,1})}{\hat{a}_{n,i^{*}}}>nK$.
But by Lemma VI.5 the LHS of this inequality is the lower bound of the sum of
the delays in the chain
$\sum^{n}_{j=2}\left(\tau^{j,l}(t_{j,l})-\tau^{j,r}(t_{j,r})\right)$. By
substitution,
$\sum^{n}_{j=2}\left(\tau^{j,l}(t_{j,l})-\tau^{j,r}(t_{j,r})\right)>nK$. It
follows that for some malicious node $j\in\\{2,\ldots,n\\}$,
$\tau^{j,l}(t_{j,l})-\tau^{j,r}(t_{j,r})>K$ which violates the delay bound
condition. ∎
Now we can show that neighbor and network discovery phases together allow the
good nodes to form a rudimentary network, where the good nodes have the same
topological view and consistent estimates of a reference clock. The first
obstacle is that the protocol is composed of stages that must be completed
sequentially by all the nodes in the network, even prior to clock
synchronization. Suppose that $[t_{k},t_{k+1})$ is the interval allocated to
the $kth$ stage. Any messages transmitted between adjacent good nodes must
arrive in the same interval they were transmitted. Since send-times are
measured with respect to the source clock, and receive-times with respect to
the destination clock, the intervals must be chosen large enough to compensate
for the maximum clock divergence caused by skew $a_{ij}\leq a_{max}$ and
offset $b_{ij}\leq a_{max}U_{0}$.
###### Lemma VI.6
There exists a sequence of adjacent time-intervals $[t_{k},t_{k+1})$ and
corresponding schedule that guarantees any message of size $W$ transmitted
(via OMC) by node $i$ in the interval $[t_{k},t_{k+1})$ (as measured by $i$’s
clock) will be received by node $j$ in the same interval as measured by node
$j$’s clock.
###### Proof:
Set $t_{k+1}:=(a_{max})^{2}t_{k}+2(a_{max})^{3}U_{0}\\\
+(a_{max})^{3}T_{MAC}(W)$. Suppose a message from node $i$ to node $j$ during
$[t_{k},t_{k+1})$ is transmitted (via the OMC) at
$t_{s}:=a_{max}t_{k}+(a_{max})^{2}U_{0}$ with respect to node $i$’s clock. By
substitution and simplification it follows that $\tau^{j}_{i}(t_{s})\geq
t_{k}$ and $\tau^{j}_{i}(t_{s}+T_{MAC}(W))<t_{k+1}$. Hence
$\tau^{j}_{i}([t_{s},t_{s}+T_{MAC}(W)))\subset[t_{k},t_{k+1})$, and so $j$
receives this message during the same interval with respect to $j$’s clock. ∎
###### Theorem VI.1
After Network Discovery, the good nodes have a common view of the topology and
consistent estimates (to within $\epsilon_{a}$) of the skew of the reference
clock.
###### Proof:
From Lemma VI.6 all good nodes will proceed through each stage of Neighbor and
Network Discovery Phases together, and therefore establish link certificates
with their good neighbors. Since they form a connected component, the good
nodes obtain a common view of their link certificates using the EIGByzMAC
algorithm and the schedule in Lemma VI.6. The good nodes can therefore infer
the network topology and the relative skews of all adjacent nodes based upon
the collection of link certificates. Using Consistency Check, the good nodes
can eliminate paths along which bad nodes have provided false skew data. The
good nodes can disseminate this information to each other using the EIGByzMAC
algorithm and Lemma VI.6 and thus obtain consistent estimates of the reference
clock to within $\epsilon_{a}$. ∎
###### Lemma VI.7
The sequence of adjacent intervals $[t_{j},t_{j+1})$, $j=0,\ldots,k$ is
contained in $[t_{0},c_{1}t_{0}+c_{2}W)$ where constants $c_{1}$ and $c_{2}$
depend on $a_{max}$, $k$, $U_{0}$, and $n$.
###### Proof:
For the OMC $T_{MAC}(W)\leq cW$, where $c$ depends on $a_{max}$, and $n$. The
result for $k=1$ follows from definition of $t_{k}$, and substitution of $cW$
into $T_{MAC}(W)$, and for general $k$ by induction and definition of $t_{k}$.
∎
###### Lemma VI.8
The time to complete Neighbor and Network Discovery Phases $T_{nei}+T_{net}$
is less than $c_{1}\log T_{life}+\frac{c_{2}}{\epsilon_{a}}$ where
$c_{1},c_{2}$ depend only on $n,a_{max},U_{0}$.
###### Proof:
From Algorithms 1, 2, 3 and 4 there are at most $6+n+n|C|+n$ protocol stages
in the Neighbor and Network Discovery Phases. Hence the time required is at
most $c_{1}t_{0}+c_{2}W$, where $W$ is the size of a message to be
transmitted, and $c_{1},c_{2}$ are constants depending on the number of
protocol stages $a_{max},U_{0},n$. The maximum size of a message is
proportional to the timing packet size $\log T_{life}$. To account for the
effect of the minimum start-time $T_{s}$ for the consistency check, we can
assume the worst case that the $T_{s}$ comes into effect during the first
protocol stage (instead of later in the Network Discovery Phase). From Theorem
V.1 the consistency check start-time is at most $\frac{c}{\epsilon_{a}}$,
where $c$ depends on $U_{0},a_{max},n$. Substitution into $t_{0}$ proves the
lemma. ∎
###### Lemma VI.9
The time required for the Data Transfer Phase is at most $c_{3}B+c_{4}D$ where
$B$ is the time spent transmitting data packets, $D$ is the size of the dead-
time separating time slots, and $c_{3},c_{4}$ depend on $n$ alone.
###### Proof:
The total number of time-slots for data transfer between all source-
destination pairs is $n^{2}(n-1)$, each supporting data transfer of size
$B_{s}$ and a dead-time $D$. ∎
###### Lemma VI.10
The time required for the Verification Phase is at most $c_{5}D$ where $c_{5}$
depends on $n$ alone.
###### Proof:
In each stage of the EIG Byzantine General’s algorithm, there are at most $n!$
vertex values that must be transmitted with each node in the neighborhood. The
value of a vertex is a list of CTVs. There are at most $2^{n}$ CTVs and at
most $n$ nodes in a CTV. Therefore the size of any message to be transmitted
by a node during EIG algorithm is at most $cD$, where $c$ is a constant
dependent on $n$. Since there are $n(n-1)$ possible source-destination pairs,
there are at most $n(n-1)$ time slots in each stage, separated at the
beginning and end by a dead-time $D$. Therefore the duration of each stage is
at most $cD+n(n-1)2D$. There are at most $n$ stages. ∎
We can now prove the main theorem of this paper.
###### Theorem VI.2
The protocol ensures that the network proceeds from startup to a functioning
network carrying data. There exists a selection of parameters $n_{iter}$, $D$,
$B$, $\epsilon_{a}$ and $T_{life}$ that achieves min-max utility over the
enabled set, to within a factor $\epsilon$, where the min is over all policies
of the bad nodes that can only adopt two actions in each CTV: conform to the
protocol and/or jam. The achieved utility is $\epsilon$-optimal.
###### Proof:
We begin by choosing parameters so that the protocol overhead, which includes
Neighbor Discovery, Network Discovery, Verification, all dead-times, and
iterations converging to the final rate vector, is an arbitrarily small
fraction of the total operating lifetime. With
$\hat{\tau}^{r}_{i}(t):=\hat{a}_{ri}t$ the estimate of reference clock $r$
with respect to the local clock at node $i$, the maximum difference in nodal
estimates is bounded as
$|\hat{\tau}^{r}_{i}(\tau^{i}(t))-\hat{\tau}^{r}_{k}(\tau^{k}_{i}(\tau^{i}(t)))|\leq
2(a_{max})^{2}\epsilon_{a}T_{life}+(a_{max})^{2}U_{0}$. With $k_{r}$ be the
number of rate vectors in the rate region, we can choose $n_{iter}$, $D$, $B$,
$\epsilon_{a}$ and $T_{life}$ to satisfy:
$\frac{n_{iter}}{n_{iter}+2^{n}k_{r}}\geq 1-\epsilon_{l}$, $\frac{B}{c_{1}\log
T_{life}+\frac{c_{2}}{\epsilon_{a}}+B+c_{3}D+c_{4}D}\geq 1-\epsilon_{d}$,
$n_{iter}((c_{1}\log T_{life}+\frac{c_{2}}{\epsilon_{a}}+B+c_{3}D+c_{4}D)\leq
T_{life}$, $2(a_{max})^{2}\epsilon_{a}T_{life}+(a_{max})^{2}U_{0}\leq D$.
These ensure that the rate loss due to failed CTVs is arbitrarily small, the
time spent transmitting data is an arbitrarily large fraction of the duration
of that iteration, the operating lifetime is large enough to support
$n_{iter}$ protocol iterations, and the dead-time $D$ is large enough to
tolerate the maximum divergence in clock estimates caused by skew error
$\epsilon_{a}$.
Let $\\{{\cal{D}}(t)\\}$ be the decreasing sequence of sets of disabled CTVs,
with limit $\bar{\cal{D}}$ attained at some finite time $T$. Suppose $x$
achieves the maximum utility for $\bar{\cal{D}}$ over the nodes in the same
component as the good nodes. No protocol can do better when $\bar{\cal{D}}$ is
disabled. The proposed protocol attains $x(1-\epsilon_{d})(1-\epsilon_{l})$. ∎
## VII Concluding Remarks
We have presented a complete suite of protocols that enables a collection of
good nodes interspersed with bad nodes to form a functioning network from
start-up, operating at a utility-optimal rate vector, regardless of what the
bad nodes conspire to do, under a certain system model. Further, the attackers
cannot decrease the utility any more than they could by just conforming to the
protocol or jamming on each CTV.
This paper is only an initial attempt to obtain a theoretical foundation for a
much needed holistic all-layer approach to secure wireless networking, and
there are several open issues. An important potential generalization is to
allow probabilistic communication. Since the protocol presented has poor
transient behavior, though overall optimal, it needs to be explored how to
increase efficiency in the transient phase.
Much further work remains to be done.
## References
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* [15] J. P. Monks, V. Bharghavan, and W.-M. Hwu. A power controlled multiple access protocol for wireless packet networks. In INFOCOM 2001. Twentieth Annual Joint Conference of the IEEE Computer and Communications Societies. Proceedings. IEEE, volume 1, pages 219–228. IEEE, 2001.
* [16] C. Perkins, E. Belding-Royer, and S. Das. Ad hoc on demand distance vector (aodv) routing. RFC 3561, July 2003.
* [17] C. E. Perkins and P. Bhagwat. Highly dynamic destination-sequenced distance-vector routing. In SIGCOMM, pages 234–244, London, UK, Aug. 1994.
* [18] J. Ponniah, Y.-C. Hu, and P. R. Kumar. An orthogonal multiple access coding scheme. Communications in Information and Systems, 12:41–76, 2012.
* [19] M. Poturalski, P. Papadimitratos, and J.-P. Hubaux. Secure neighbor discovery in wireless networks: formal investigation of possibility. ASIACCS ’08, pages 189–200, New York, NY, USA, 2008. ACM.
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|
arxiv-papers
| 2013-09-11T17:47:57 |
2024-09-04T02:49:50.745450
|
{
"license": "Public Domain",
"authors": "Jonathan Ponniah, Yih-Chun Hu and P. R. Kumar",
"submitter": "P. R. Kumar",
"url": "https://arxiv.org/abs/1309.2904"
}
|
1309.2963
|
∎
11institutetext: P. Shakarian 22institutetext: Network Science Center and
Dept. Electrical Engineering and Computer Science
U.S. Military Academy
West Point, NY 10996
Tel.: 845-938-5576
22email: [email protected] 33institutetext: S. Eyre 44institutetext: Network
Science Center and
Dept. Electrical Engineering and Computer Science
U.S. Military Academy
West Point, NY 10996
Tel.: 845-938-5576
44email: [email protected] 55institutetext: Damon Paulo
Network Science Center and
Dept. Electrical Engineering and Computer Science
U.S. Military Academy
West Point, NY 10996
Tel.: 845-938-5576
55email: [email protected]
# A Scalable Heuristic for Viral Marketing Under the Tipping Model
Paulo Shakarian Sean Eyre Damon Paulo
###### Abstract
In a “tipping” model, each node in a social network, representing an
individual, adopts a property or behavior if a certain number of his incoming
neighbors currently exhibit the same. In viral marketing, a key problem is to
select an initial ”seed” set from the network such that the entire network
adopts any behavior given to the seed. Here we introduce a method for quickly
finding seed sets that scales to very large networks. Our approach finds a set
of nodes that guarantees spreading to the entire network under the tipping
model. After experimentally evaluating 31 real-world networks, we found that
our approach often finds seed sets that are several orders of magnitude
smaller than the population size and outperform nodal centrality measures in
most cases. In addition, our approach scales well - on a Friendster social
network consisting of 5.6 million nodes and 28 million edges we found a seed
set in under 3.6 hours. Our experiments also indicate that our algorithm
provides small seed sets even if high-degree nodes are removed. Lastly, we
find that highly clustered local neighborhoods, together with dense network-
wide community structures, suppress a trend’s ability to spread under the
tipping model.
###### Keywords:
social networks viral marketing tipping model
## 1 Introduction
A much studied model in network science, tippingGran78 ; Schelling78 ; jy05
(a.k.a. deterministic linear thresholdkleinberg ) is often associated with
“seed” or “target” set selection, chen09siam (a.k.a. the maximum influence
problem). In this problem, we have a social network in the form of a directed
graph and thresholds for each individual. Based on this data, the desired
output is the smallest possible set of individuals (seed set) such that, if
initially activated, the entire population will become activated (adopting the
new property). This problem is NP-Complete kleinberg ; Dreyer09 so
approximation algorithms must be used. Though some such algorithms have been
proposed, leskovec07 ; chen09siam ; benzwi09 ; chen10 none seem to scale to
very large data sets. Here, inspired by shell decomposition, ShaiCarmi07032007
; InfluentialSpreaders_2010 ; baxter11 we present a method guaranteed to find
a set of nodes that causes the entire population to activate - but is not
necessarily of minimal size. We then evaluate the algorithm on $31$ large,
real-world, social networks and show that it often finds very small seed sets
(often several orders of magnitude smaller than the population size). We also
show that the size of a seed set is related to Louvain modularity and average
clustering coefficient. Therefore, we find that dense community structure
combined with tight-knit local neighborhoods inhibit the spreading of
activation under the tipping model. We also found that our algorithm
outperforms the classic centrality measures and is robust against the removal
of high-degree nodes.
The rest of the paper is organized as follows. In Section 2, we provide formal
definitions of the tipping model. This is followed by the presentation of our
new algorithm in Section 3. We then describe our experimental results in
Section 4. Finally, we provide an overview of related work in Section 5.
## 2 Technical Preliminaries
Throughout this paper we assume the existence of a social network, $G=(V,E)$,
where $V$ is a set of vertices and $E$ is a set of directed edges. We will use
the notation $n$ and $m$ for the cardinality of $V$ and $E$ respectively. For
a given node $v_{i}\in V$, the set of incoming neighbors is $\eta^{in}_{i}$,
and the set of outgoing neighbors is $\eta^{out}_{i}$. The cardinalities of
these sets (and hence the in- and out-degrees of node $v_{i}$) are
$d^{in}_{i},d^{out}_{i}$ respectively. We now define a threshold function that
for each node returns the fraction of incoming neighbors that must be
activated for it to become activate as well.
###### Definition 1 (Threshold Function)
We define the threshold function as mapping from V to $(0,1]$. Formally:
$\theta:V\rightarrow(0,1]$.
For the number of neighbors that must be active, we will use the shorthand
$k_{i}$. Hence, for each $v_{i}$, $k_{i}=\lceil\theta(v_{i})\cdot
d^{in}_{i}\rceil$. We now define an activation function that, given an initial
set of active nodes, returns a set of active nodes after one time step.
###### Definition 2 (Activation Function)
Given a threshold function, $\theta$, an activation function $A_{\theta}$ maps
subsets of V to subsets of V, where for some $V^{\prime}\subseteq V$,
$A_{\theta}(V^{\prime})=V^{\prime}\cup\\{v_{i}\in V\ s.t.\ |\eta^{in}_{i}\cap
V^{\prime}|\geq k_{i}\\}$ (1)
We now define multiple applications of the activation function.
###### Definition 3 (Multiple Applications of the Activation Function)
Given a natural number $i>0$, set $V^{\prime}\subseteq V$, and threshold
function, $\theta$, we define the multiple applications of the activation
function, ${A^{i}_{\theta}}(V^{\prime})$, as follows:
$A^{i}_{\theta}(V^{\prime})=\begin{cases}A_{\theta}(V^{\prime})&\text{if
$i=1$}\\\
A_{\theta}(A^{i-1}_{\theta}(V^{\prime}))&\text{otherwise}\end{cases}$ (2)
Clearly, when $A^{i}_{\theta}(V^{\prime})=A^{i-1}_{\theta}(V^{\prime})$ the
process has converged. Further, this always converges in no more than $n$
steps (as, prior to converging, a process must, in each step, activate at
least one new node). Based on this idea, we define the function $\Gamma$ which
returns the set of all nodes activated upon the convergence of the activation
function.
###### Definition 4 ($\Gamma$ Function)
Let j be the least value such that
$A^{j}_{\theta}(V^{\prime})=A^{j-1}_{\theta}(V^{\prime})$. We define the
function $\Gamma_{\theta}:2^{V}\rightarrow 2^{V}$ as follows.
$\mathbf{\Gamma_{\theta}}(V^{\prime})=A^{j}_{\theta}(V^{\prime})$ (3)
We now have all the pieces to introduce our problem - finding the minimal
number of nodes that are initially active to ensure that the entire set $V$
becomes active.
###### Definition 5 (The MIN-SEED Problem)
The MIN-SEED Problem is defined as follows: given a threshold function,
$\theta$, return $V^{\prime}\subseteq V\ s.t.\ \Gamma_{\theta}(V^{\prime})=V$,
and there does not exist $V^{\prime\prime}\subseteq V$ where
$|V^{\prime\prime}|<|V^{\prime}|$ and $\Gamma_{\theta}(V^{\prime\prime})=V$.
The following theorem is from the literature kleinberg ; Dreyer09 and tells
us that the MIN-SEED problem is NP-complete.
###### Theorem 2.1 (Complexity of MIN-SEED kleinberg ; Dreyer09 )
MIN-SEED in NP-Complete.
## 3 Algorithms
In this section, we introduce an integer program that solved the MIN-SEED
problem exactly and our new decomposition-based heuristic.
### 3.1 Exact Approach
Below we present SEED-IP, an integer program that if solved exactly,
guarantees an exact solution to MIN-SEED (see Proposition 1). Though, in
general, solving an integer program is also NP-hard, suggesting that an exact
solution will likely take exponential time, good approximation techniques such
as branch-and-bound exist and mature tools such as QSopt and CPLEX can readily
take and approximate solutions to integer programs.
###### Definition 6 (SEED-IP)
$\displaystyle\min\sum_{i}x_{i,1},\,\,\,\,\,\,\,\,\,\,\mathit{w.r.t.}$ (4)
$\displaystyle\forall i,t\in\\{1,\ldots,n\\},$ $\displaystyle
x_{i,t}\in\\{0,1\\}$ (5) $\displaystyle\forall i,$ $\displaystyle x_{i,n}=1$
(6) $\displaystyle\forall i,\forall t>0,$
$\displaystyle\,\,\,\,\,\,\,\,\,\,x_{i,t}\leq
x_{i,t-1}+\frac{1}{d^{in}_{i}\theta(v_{i})}\sum_{v_{j}\in\eta^{in}_{i}}x_{j,t-1}$
(7)
###### Proposition 1
If $V^{\prime}$ is a solution to MIN-SEED, then setting $\forall v_{i}\in
V^{\prime},x_{i,1}=1$ and $\forall v_{i}\notin V^{\prime},x_{i,1}=0$ is a
solution to SEED-IP.
If the vector $[x_{i,t}]$ is a solution to SEED-IP, then
$\\{v_{i}|x_{i,1}=1\\}$ is a solution to MIN-SEED.
###### Proof
Claim 1: If $V^{\prime}$ is a solution to MIN-SEED, then setting $\forall
v_{i}\in V^{\prime},x_{i,1}=1$ and $\forall v_{i}\notin V^{\prime},x_{i,1}=0$
is a solution to SEED-IP.
Let $[x_{i,t}]$ be a vector for SEED-IP created as per claim 1. Suppose, by
way of contradiction (BWOC), there exists vector $[x^{\prime}_{i,t}]$ s.t.
$\sum_{i}x^{\prime}_{i,1}<\sum_{i}x_{i,1}$. However, consider the set of nodes
$V^{\prime\prime}=\\{v_{i}|x^{\prime}_{i,1}=1\\}$. By Constraint 7 of SEED-IP,
we know that, for $t>1$, that if $x^{\prime}_{i,t}=1$, we have $v_{i}\in
A^{t}_{\theta}(V^{\prime\prime})$. Hence, by Constaint 6 $V^{\prime\prime}$ is
a solution to MIN-SEED. This means that $|V^{\prime\prime}|<|V^{\prime}|$ as
$\sum_{i}x^{\prime}_{i,1}<\sum_{i}x_{i,1}$, which is a contradiction.
Claim 2: If the vector $[x_{i,t}]$ is a solution to SEED-IP, then
$\\{v_{i}|x_{i,1}=1\\}$ is a solution to MIN-SEED.
Suppose, BWOC, there exists set $V^{\prime\prime}$ that is a solution to MIN-
SEED s.t. $|V^{\prime\prime}|<|\\{v_{i}|x_{i,1}=1\\}|$. Consider the vector
$[x^{\prime}_{i,t}]$ where $\forall i,$ $x^{\prime}_{i,0}=1$ iff $v_{i}\in
V^{\prime\prime}$. By Constraint 7 of SEED-IP, we know that, for $t>1$, that
if $v_{i}\in A^{t}_{\theta}(V^{\prime\prime})$, we have $x^{\prime}_{i,t}=1$.
Hence, as $A^{t}_{\theta}(V^{\prime\prime})=V$, know that $[x^{\prime}_{i,t}]$
satisfies Constraint 6. Hence, as
$|V^{\prime\prime}|<|\\{v_{i}|x_{i,1}=1\\}|$, we know
$\sum_{i}x^{\prime}_{i,1}<\sum_{i}x_{i,1}$, which is a contradiction.
Proof of theorem: Follows directly form claims 1-2.
However, despite the availability of approximate solvers, SEED-IP requires a
quadratic number of variables and constraints (Proposition 2), which likely
will prevent this approach from scaling to very large datasets. As a result,
in the next section we introduce our heuristic approach.
###### Proposition 2
SEED-IP requires $n^{2}$ variables and $2n^{2}$ constraints.
### 3.2 Heuristic
To deal with the intractability of the MIN-SEED problem, we design an
algorithm that finds a non-trivial subset of nodes that causes the entire
graph to activate, but we do not guarantee that the resulting set will be of
minimal size. The algorithm is based on the idea of shell decomposition often
cited in physics literature Seidman83 ; ShaiCarmi07032007 ;
InfluentialSpreaders_2010 ; baxter11 but modified to ensure that the
resulting set will lead to all nodes being activated. The algorithm,
TIP_DECOMP is presented in this section.
Algorithm 1 TIP_DECOMP
0: Threshold function, $\theta$ and directed social network $G=(V,E)$
0: $V^{\prime}$
1: For each vertex $v_{i}$, compute $k_{i}$.
2: For each vertex $v_{i},\ dist_{i}=d_{i}^{in}-k_{i}$.
3: FLAG = TRUE.
4: while FLAG do
5: Let $v_{i}$ be the element of $v$ where $dist_{i}$ is minimal.
6: if $dist_{i}=\infty$ then
7: FLAG = FALSE.
8: else
9: Remove $v_{i}$ from $G$ and for each $v_{j}$ in $\eta_{i}^{out}$, if
$dist_{j}>0$, set $dist_{j}=dist_{j}-1$. Otherwise set $dist_{j}=\infty$.
10: end if
11: end while
12: return All nodes left in $G$.
Intuitively, the algorithm proceeds as follows (Figure 1). Given network
$G=(V,E)$ where each node $v_{i}$ has threshold
$k_{i}=\lceil\theta(v_{i})\cdot d^{in}_{i}\rceil$, at each iteration, pick the
node for which $d^{in}_{i}-k_{i}$ is the least but positive (or $0$) and
remove it. Once there are no nodes for which $d^{in}_{i}-k_{i}$ is positive
(or $0$), the algorithm outputs the remaining nodes in the network.
Figure 1: Example of our algorithm for a simple network depicted in box A. We
use a threshold value set to $50\%$ of the node degree. Next to each node
label (lower-case letter) is the value for $d^{in}_{i}-k_{i}$ (where
$k_{i}=\lceil\frac{d^{in}_{i}}{2}\rceil$). In the first four iterations, nodes
e, f, h, and i are removed resulting in the network in box B. This is followed
by the removal of node j resulting in the network in box C. In the next two
iterations, nodes a and b are removed (boxes D-E respectively). Finally, node
c is removed (box F). The nodes of the final network, consisting of d and g,
have negetive values for $d_{i}-\theta_{i}$ and become the output of the
algorithm.
Now, we prove that the resulting set of nodes is guaranteed to cause all nodes
in the graph to activate under the tipping model. This proof follows from the
fact that any node removed is activated by the remaining nodes in the network.
###### Theorem 3.1
If all nodes in $V^{\prime}\ \subseteq\ V$ returned by TIP_DECOMP are
initially active, then every node in $V$ will eventually be activated, too.
###### Proof
Let $w$ be the total number of nodes removed by TIP_DECOMP, where $v_{1}$ is
the last node removed and $v_{w}$ is the first node removed. We prove the
theorem by induction on $w$ as follows. We use $P(w)$ to denote the inductive
hypothesis which states that all nodes from $v_{1}$ to $v_{w}$ are active. In
the base case, $P(1)$ trivially holds as we are guaranteed that from set
$V^{\prime}$ there are at least $k_{1}$ edges to $v_{1}$ (or it would not be
removed). For the inductive step, assuming $P(w)$ is true, when $v_{w+1}$ was
removed from the graph $dist_{w+1}\geq 0$ which means that $d_{w+1}^{in}\geq
k_{w+1}$. All nodes in $\eta^{in}_{w+1}$ at the time when $v_{w+1}$ was
removed are now active, so $v_{w+1}$ will now be activated - which completes
the proof.
We also note that by using the appropriate data structure (we used a binomial
heap in our implementation), for a network of $n$ nodes and $m$ edges, this
algorithm can run in time $O(m\log n)$.
###### Proposition 3
The complexity of TIP_DECOMP is $O(m\cdot log(n))$.
###### Proof
If we use a binomial heap as described in cormen , we can create a heap where
we store each node and assign it a key value of $dist_{i}$ for each node
$v_{i}$. The creation of a heap takes constant time and inserting the $n$
vertices will take $O(nlog(n))$ time. We can also maintain a list data
structure as well. In the course of the while loop, all nodes will either be
removed (as per the algorithm), decreased in key-value no more than
$d^{in}_{i}$ or increased to infinity (which we can implement as being removed
and added to the list). Hence, the number of decrease key or removal
operations is bounded by $n+\sum_{i}d^{in}_{i}$. As $\sum_{i}d^{in}_{i}=m$
(where $m$ is the number of edges). As $O(m\cdot log(n))$, the statement
follows.
## 4 Results
In this section we describe the results of our experimental evaluation. We
describe the datasets we used for the experiments in Section 4.1. We evaluate
the run-time of TIP_DECOMP in Section 4.1.5. In Section 4.1.6, we evaluate the
size of the seed-set returned by the algorithm and we compare this to the seed
size returned by known centrality measures in Section 4.2. The speed of the
activiation process initiated with seed sets discovered by our algorithm is
described in Section 4.3. We then study how the removal of high-degree nodes
and community structure affect the results of the algorithm in Sections 4.4
and 4.4.1 respectively.
The algorithm TIP_DECOMP was written using Python 2.6.6 in 200 lines of code
that leveraged the NetworkX library available from
http://networkx.lanl.gov/. The code used a binomial heap library written by
Björn B. Brandenburg available from http://www.cs.unc.edu/$\sim$bbb/. The
experiments were run on a computer equipped with an Intel X5677 Xeon Processor
operating at 3.46 GHz with a 12 MB Cache running Red Hat Enterprise Linux
version 6.1 and equipped with 70 GB of physical memory. All statistics
presented in this section were calculated using R 2.13.1.
### 4.1 Datasets
In total, we examined $36$ networks: nine academic collaboration networks,
three e-mail networks, and $24$ networks extracted from social-media sites.
The sites included included general-purpose social-media (similar to Facebook
or MySpace) as well as special-purpose sites (i.e. focused on sharing of
blogs, photos, or video).
All datasets used in this paper were obtained from one of four sources: the
ASU Social Computing Data Repository, Zafarani+Liu:2009 the Stanford Network
Analysis Project, snap the University of Michigan, umich and Universitat
Rovira i Virgili.uvi $31$ of the networks considered were symmetric – i.e. if
a directed edge from vertex $v$ to $v^{\prime}$ exists, there is also an edge
from vertex $v^{\prime}$ to $v$. Tables 1 (A-C) show some of the pertinent
qualities of the symmetric networks. The networks are categorized by the
results for the MIN-SEED experiments (explained later in this section).
Additionally, we also looked at several non-symmetric (directed) networks and
placed them in their own category. In what follows, we provide their real-
world context.
#### 4.1.1 Category A
* •
BlogCatalog is a social blog directory that allows users to share blogs with
friends. Zafarani+Liu:2009 The first two samples of this site, BlogCatalog1
and 2, were taken in Jul. 2009 and June 2010 respectively. The third sample,
BlogCatalog3 was uploaded to ASU’s Social Computing Data Repository in Aug.
2010.
* •
Buzznet is a social media network designed for sharing photographs, journals,
and videos. Zafarani+Liu:2009 It was extracted in Nov. 2010.
* •
Douban is a Chinese social medial website designed to provide user reviews
and recommendations. Zafarani+Liu:2009 It was extracted in Dec. 2010.
* •
Flickr is a social media website that allows users to share photographs.
Zafarani+Liu:2009 It was uploaded to ASU’s Social Computing Data Repository
in Aug. 2010.
* •
Flixster is a social media website that allows users to share reviews and
other information about cinema. Zafarani+Liu:2009 It was extracted in Dec.
2010.
* •
FourSquare is a location-based social media site. Zafarani+Liu:2009 It was
extracted in Dec. 2010.
* •
Frienster is a general-purpose social-networking site. Zafarani+Liu:2009 It
was extracted in Nov. 2010.
* •
Last.Fm is a music-centered social media site. Zafarani+Liu:2009 It was
extracted in Dec. 2010.
* •
LiveJournal is a site designed to allow users to share their blogs.
Zafarani+Liu:2009 It was extracted in Jul. 2010.
* •
Livemocha is touted as the “world’s largest language community.”
Zafarani+Liu:2009 It was extracted in Dec. 2010.
* •
WikiTalk is a network of individuals who set and received messages while
editing WikiPedia pages. snap It was extracted in Jan. 2008.
#### 4.1.2 Category B
* •
Delicious is a social bookmarking site, designed to allow users to share web
bookmarks with their friends. Zafarani+Liu:2009 It was extracted in Dec.
2010.
* •
Digg is a social news website that allows users to share stories with friends.
Zafarani+Liu:2009 It was extracted in Dec. 2010.
* •
EU E-Mail is an e-mail network extracted from a large European Union research
institution. snap It is based on e-mail traffic from Oct. 2003 to May 2005.
* •
Hyves is a popular general-purpose Dutch social networking site.
Zafarani+Liu:2009 It was extracted in Dec. 2010.
* •
Yelp is a social networking site that allows users to share product reviews.
Zafarani+Liu:2009 It was extracted in Nov. 2010.
#### 4.1.3 Category C
* •
CA-AstroPh is a an academic collaboration network for Astro Physics from Jan.
1993 - Apr. 2003. snap
* •
CA-CondMat is an academic collaboration network for Condense Matter Physics.
Samples from 1999 (CondMat99), 2003 (CondMat03), and 2005 (CondMat05) were
obtained from the University of Michigan. umich A second sample from 2003
(CondMat03a) was obtained from Stanford University. snap
* •
CA-GrQc is a an academic collaboration network for General Relativity and
Quantum Cosmology from Jan. 1993 - Apr. 2003. snap
* •
CA-HepPh is a an academic collaboration network for High Energy Physics -
Phenomenology from Jan. 1993 - Apr. 2003. snap
* •
CA-HepTh is a an academic collaboration network for High Energy Physics -
Theory from Jan. 1993 - Apr. 2003. snap
* •
CA-NetSci is a an academic collaboration network for Network Science from May
2006.
* •
Enron E-Mail is an e-mail network from the Enron corporation made public by
the Federal Energy Regulatory Commission during its investigation. snap
* •
URV E-Mail is an e-mail network based on communications of members of the
University Rovira i Virgili (Tarragona). uvi It was extracted in 2003.
* •
YouTube is a video-sharing website that allows users to establish friendship
links. Zafarani+Liu:2009 The first sample (YouTube1) was extracted in Dec.
2008. The second sample (YouTube2) was uploaded to ASU’s Social Computing Data
Repository in Aug. 2010.
#### 4.1.4 Non-Symmetric Networks
* •
Epinions is a consumer review website that allows members to establish
directed trust relationships. snap
* •
WikiVote is a sample of Wikipedia users voting beahavior (who votes for whom).
snap
* •
Slashdot formerly had a feature called “Slashdot Zoo” that allowed users to
tag each other as friend or foe. We looked at three samples based on
friendship relationships: one sample from 2008 (Slashdot1) and two from 2009
(Slashdot2-Slashdot3). snap
Table 1: Information on the networks in Categories A, B, and C. Table 2:
Information on non-symmetric networks.
#### 4.1.5 Runtime
First, we examined the runtime of the algorithm (see Figure 2 and Table 3).
Our experiments aligned well with our time complexity result (Proposition 3).
For example, a network extracted from the Dutch social-media site Hyves
consisting of $1.4$ million nodes and $5.5$ million directed edges was
processed by our algorithm in at most $12.2$ minutes. The often-cited
LiveJournal dataset consisting of $2.2$ million nodes and $25.6$ million
directed edges was processed in no more than $66$ minutes - a short time to
approximate an NP-hard combinatorial problem on a large-sized input.
Figure 2: $m\ln n$ vs. runtime in seconds (log scale, $m$ is number of edges,
$n$ is number of nodes). The relationship is linear with $R^{2}=0.9015$,
$p=2.2\cdot 10^{-16}$. Table 3: Runtime data on the datasets used in the
experiments.
#### 4.1.6 Seed Size
For each network, we performed $10$ “integer” trials. In these trials, we set
$\theta(v_{i})=\min(d^{in}_{i},k)$ where $k$ was kept constant among all
vertices for each trial and set at an integer in the interval $[1,10]$. We
evaluated the ability of a network to promote spreading under the tipping
model based on the size of the set of nodes returned by our algorithm (as a
percentage of total nodes). For purposes of discussion, we have grouped our
networks into three categories based on results (Figure 3 and Table 4). We
have also included results for symmetric networks (Figure 4 and Table 5). In
general, online social networks had the smallest seed sets - $13$ networks of
this type had an average seed set size less than $2\%$ of the population
(these networks were all in Category A). We also noticed, that for most
networks, there was a linear realtion between threshold value and seed size.
Figure 3: Threshold value (assigned as an integer in the interval $[1,10]$)
vs. size of initial seed set as returned by our algorithm in our three
identified categories of networks (categories A-C are depicted in panels A-C
respectively). Average seed sizes were under $2\%$ for Categorty A, $2-10\%$
for Category B and over $10\%$ for Category C. The relationship, in general,
was linear for categories A and B and lograthimic for C. CA-NetSci had the
largest Louvain Modularity and clustering coefficient of all the networks.
This likely explains why that particular network seems to inhibit spreading.
Figure 4: Threshold value assigned as both an integer in the range $[1,10]$
(panel 1) as well as a fraction of node degree (panel 2) for the non-symmetric
networks.
Category A can be thought of as social networks highly susceptible to
influence - as a very small fraction of initially activated individuals can
lead to activation of the entire population. All were extracted from social
media websites. For some of the lower threshold levels, the size of the set of
seed nodes was particularly small. For a threshold of three, $11$ of the
Category A networks produced seeds smaller than $0.5\%$ of the total
populationa. For a threshold of four, nine networks met this criteria.
Networks in Category B are susceptible to influence with a relatively small
set of initial nodes - but not to the extent of those in Category A. They had
an average initial seed size greater than $2\%$ but less than $10\%$. Members
in this group included two general purpose social media networks, two
specialty social media networks, and an e-mail network. Non-symmetric networks
generally perofrmed somewhat poorer than Category B networks (though in
general, not as poorly as those in Category C). The initial seed sizes for the
non-symmmetric networks ranged from $3\%$ to $29\%$.
Category C consisted of networks that seemed to hamper diffusion in the
tipping model, having an average initial seed size greater than $10\%$. This
category included all of the academic collaboration networks, two of the email
networks, and two networks derived from friendship links on YouTube.
We also studied the effects on spreading when the threshold values were
assigned as a specific fraction of each node’s in-degree jy05 ; wattsDodds07 ,
which results in heterogeneous $\theta_{i}$’s across the network. We performed
$12$ trials for each network. Thresholds for each trial were based on the
product of in-degree and a fraction in the interval $[0.05,0.60]$ (multiples
of $0.05$). The results (Figure 5 and Table 4; for non-symmertic networks see
Figure 4 and Table 5) were analogous to our integer tests. We also compared
the averages over these trials with $M$ and $C$ and obtained similar results
as with the other trials (Figure 14 B).
Figure 5: Threshold value (assigned as a fraction of node in-degree as a
multiple of $0.05$ in the interval $[0.05,0.60]$) vs. size of initial seed set
as returned by our algorithm in our three identified categories of networks
(categories A-C are depicted in panels A-C respectively, categories are the
same as in Figure 1). Average seed sizes were under $5\%$ for Categorty A,
$1-7\%$ for Category B and over $3\%$ for Category C. In general, the
relationship between threshold and initial seed size for networks in all
categories was exponential.
### 4.2 Comparison with Centrality Measures
We compared our results with six popular centrality measures: degree,
betweenness, closeness, shell number, eigenvector, and PageRank. Here, we
define degree centrality is simply the number of outgoing adjacent nodes.
111Note that in the symmetric networks we examined, our empirical results hold
for the number of incoming adjacent edges as well as the total number of
adjacent edges. The intuition behind high betweenness centrality nodes is that
they function as “bottlenecks” as many paths in the network pass through them.
Hence, betweenness is a medial centrality measure. Let $\sigma_{st}$ be the
number of shortest paths between nodes $s$ and $t$ and $\sigma_{st}(v)$ be the
number of shortest paths between $s$ and $t$ containing node $v$. In freeman77
, betweenness centrality for node $v$ is defined as $\sum_{s\neq v\neq
t}\frac{\sigma_{st}(v)}{\sigma_{st}}$. In most implementations, including the
ones used in this paper, the algorithm of brandes01 is used to calculate
betweenness centrality. Another common measure from the literature that we
examined is closeness freeman79cent . Given node $i$, its closeness $C_{c}(i)$
is the inverse of the average shortest path length from node $i$ to all other
nodes in the graph. Intuitively, closeness measures how “close” it is to all
other nodes in a network. Formally, if we define the shortest path between
nodes $i$ to $j$ as function $d_{G}(i,j)$, we can express the average path
length from $i$ to all other nodes as
$L_{i}=\frac{\sum_{j\in V\setminus i}d_{G}(i,j)}{|V|-1}.$ (8)
Hence, the closeness of a node can be formally written as
$C_{c}(i)=\frac{1}{L_{i}}=\frac{|V|-1}{\sum_{j\in V\setminus i}d_{G}(i,j)}.$
(9)
The idea of shell number is based on a core to which a node lies in. A
$c$-core of a network is the subgraph in which every node is connected to the
rest of the network by at least $c$ edges. A node is assigned a shell number
based on the maximal core that contains it. This value can be derived exactly
using shell decomposition Seidman83 . The eigenvector centrality bona72 of a
node is assigned based on the associated entry in the eigenvector of the
adjacency matrix corresponding to the largest real eigenvalue. The PageRank
Page98 for each node based on the PageRank of its neighbors. An initial value
for rank is considered for each node and the relationship is then computed
iteratively until convergence is reached. Intuitively, PageRank can be thought
of as the importance of a node based on the importance of its neighbors. We
note that shell number, eigenvector, and PageRank are often associated with
diffusion processes. A more complete discussion of centrality measures can be
found in wasserman1994social .
We evaluated the performance of centrality measures in finding a seed set by
iteratively selecting the most central nodes with respect to a given measure
until the $\Gamma_{\theta}$ of that set returns the set of all nodes. Due to
the computational overhead of calculating these centrality measures and the
repeated re-evaluation of $\Gamma_{\theta}$, we limited this comparison to
only BlogCatalog3, CA-HepTh, CA-NetSci, URV E-Mail, and Douban (no betweeness
calcualted for Douban). As with the experiments in the previous section, we
studied threshold settings based on an integer in the interval $[1,10]$ (see
Figure 6) and as a fraction of incoming neighbors in the interval
$[0.05,0.60]$ (multiples of $0.05$, see Figure 8). In general, selecting
highly-central nodes is an inefficient method for finding small seed sets.
In all but the lowest threshold settings, the use of centrality measures for
the integer-threshold trials proved to significantly underperformed when the
method presented in this paper - often returning seed-sets several orders of
magnitude larger and in many cases including the majority of nodes in the
network. Even for the centrality measures outperformed our method in these
trials, the reduction in seed set size was minimal (the greatest reduction in
seed set size experienced in a centrality-measure test over the algorithm of
this paper was $0.09\%$, while often producing seed sets orders of magnitude
greater than our method). This held even for the centrality measures
associated with diffusion (shell number, eigenvector, and PageRank).
Figure 6: The use of degree, betweenness and closeness to find seed-sets on
select networks when the threshold is set to an integer in the interval
$[1,10]$. For these trials, centrality measure returned significantly larger
(several orders of magnitude) larger seed sets than our approach. Figure 7:
The use of shell number, Eigenvector, and PageRank to find seed-sets on select
networks when the threshold is set to an integer in the interval $[1,10]$.
Our tests using fractional-based thresholds tell a slightly different story.
While our method still generally outperformed the centrality measures for the
fractional tests, there were a few cases where the centrality measures fared
better. With BlogCatalog3 all of the centrality measures outperformed our
algorithm in the fraction-based experiments. For that dataset, centrality-
based algorithm consistently outperformed our method finding seed sets with
less members (by $3.13-3.29\%$ of the population, on average). With URV-Email,
many trials that utilized a lower threshold setting outperformed our method,
but never finding a feed set with smaller by more than $8\%$ of the total
population. However, in the larger threshold settings, our method consistently
found smaller seeds. For a given centrality measure for this dataset,
centrality measures on average provided poorer results than our algorithm
ranged - returning seed sets which were, on average $10.22-67.14\%$ (by
overall population) larger than that returned by our algorithm. Perhaps the
most interesting result among the centrality measures were the PageRank
fraction-based tests on CA-NetSci, which is associated with the largest seed
sets. PageRank found seed sets that were, on average $14.45\%$ smaller (by
population) than our approach. Additionally, though centrality measures
outperformed TIP_DECOMP for BlogCatalog3, this does not appear to hold for all
social networks as the seed sets returned using centrality measures for the
Douban approaches at least an order of magnitude increase over our method for
nearly every fractional threshold setting for all centrality measures. Hence,
we conclude that for fraction-based thresholds, using centrality measures to
find seed sets provides inconsistent results, and when it fails, it tends to
provide a large portion of the network. A possibility for a practical
algorithm that could combine both methods would be to first run TIP_DECOMP,
returning some set $V^{\prime}$. Then, create set $V^{\prime\prime}$ by
selecting the most central nodes until either
$|V^{\prime}|=|V^{\prime\prime}|$ or $\Gamma_{\theta}(V^{\prime\prime})=V$
(whichever ensures the lower cardinality for $V^{\prime\prime}$. If
$|V^{\prime}|=|V^{\prime\prime}|$, return $V^{\prime}$, otherwise return
$V^{\prime\prime}$. For such an approach, we would likely recommend using
degree centrality due to its ease of computation and performance in our
experiments. However, we note that highly-central nodes often may not be
realistic targets for a viral-marketing campaign. For instance, it may be
cost-prohibitive to create a seed set consisting of major celebrities in order
to spread the use of a product. As such is a practical concern, we look at the
performance of TIP_DECOMP when high-degree nodes are removed in the next
section.
Figure 8: The use of degree, betweenness and closeness to find seed-sets on
select networks when the threshold is set to an fraction in the interval
$[0.05,0.60]$ (multiples of $0.05$). Figure 9: The use of shell number,
Eigenvector, and PageRank to find seed-sets on select networks when the
threshold is set to an fraction in the interval $[0.05,0.60]$ (multiples of
$0.05$).
### 4.3 The Speed of the Activation Process and Sets of “Critical Mass”
An important aspect to consider in viral marketing is the speed of the
activation process. We illustrate this speed for several networks under a
threshold of $2$ as well as a majority threshold (half of each nodes
neighbors) in Figure 10. Interestingly, we found that the size of the initial
seed set was not indicative of the speed of spreading. For instance, in
BlogCatalog3, a Category A network (for which our algorithm found a very small
seed set) the activation process proceeded quickly when compared to the others
examined. However, this was also true for CA-NetSci, a Category C network
(large seed set). Conversely, the activation process in the Douban and CA-
HepTh networks (also Category A and C, respectively) proceeded more slowly
than the rest.
Another interesting feature we learned in exploring the speed of the
activation process was that in all of our experiments there was a single time
step where the number of activated nodes increased significantly more than the
other time periods - sometimes by several orders of magnitude (see Figure 11).
We can think of such a set of activated nodes as when the population reaches a
“critical mass” which results in mass adoption in the next interval. In many
cases, such a critical mass is reached early - normally in the first few time-
steps.
Finding a subset of the population of “critical mass” may be an important
problem in its own right. Though the critical mass point will often be larger
than the seed set found by an algorithm in this paper, we can be assured that
in one time step of the model, the number of individuals reached (with a
certain number of signals from their neighbors) is substantially larger than
the investment. In practice, this could lead to quicker spreading of
information in an advertising campaign, for example. Further, our experiments
indicate that order-of-magnitude critical mass sets exist in several networks.
We are currently conducting further research on this topic.
Figure 10: An examination of several of speed of activation initiated from the
seed set using a threshold of two (panel 1) and a majority threshold (panel
2). Figure 11: Greatest Percent increase experienced in a single time step
(the effect of reaching “critical mass”) for integer-based and percentage-
based thresholds (panel 1 and 2 respectively).
### 4.4 Effect of Removing High-Degree Nodes
In the last section we noted that high-degree nodes may not always be
targetable in a viral marketing campaign (i.e. it may be cost prohibitive to
include them in a seed set). In this section, we explore the affect of
removing high-degree nodes on the size of the seed-set returned by TIP_DECOMP.
This type of node removal has also recently been studied in a different
context in nodeRemRef . In these trials, we studied all $31$ networks and
looked at two specific threshold settings: an integer threshold of $2$ (Figure
12) and a fractional threshold of $0.5$ (Figure 13). We then studied the
effect of removing up to $50\%$ of the nodes in order from greatest to least
degree.
With an integer threshold of $2$, networks in category A still retained a
seed-size (as returned by TIP_DECOMP) two orders of magnitude smaller than the
population size up to the removal of $10\%$ of the top degree nodes, and for
many networks this was maintained to $50\%$. Networks in category B retained
seed sets an order of magnitude smaller than the population for up to $50\%$
of the nodes removed. For most networks in category C, the seed size remained
about a quarter of the population size up to $15\%$ of the top degree nodes
being removed.
Figure 12: Size of the seed set returned by TIP_DECOMP (as a fraction of the
popualtion) as a function of the percent of the highest degree nodes removed
from the network with an integer theshold of $2$ for networks in categories
A-C.
With a fractional threshold of $0.5$, we noted that many networks in category
A actually had larger seed sets (as returned by TIP_DECOMP) when more high
degree nodes are removed. Further, networks in categories A-B retained seed
sets of at least an order of magnitude smaller than the population in these
tests while most networks in category C retained sizes of about a quarter of
the population.
Figure 13: Size of the seed set returned by TIP_DECOMP (as a fraction of the
popualtion) as a function of the percent of the highest degree nodes removed
from the network with an fractional theshold of $0.5$ for networks in
categories A-C.
#### 4.4.1 Seed Size as a Function of Community Structure
In this section, we view the results of our heuristic algorithm as a
measurement of how well a given network promotes spreading. Here, we use this
measurement to gain insight into which structural aspects make a network more
likely to be “tipped.” We compared our results with two network-wide measures
characterizing community structure. First, clustering coefficient ($C$) is
defined for a node as the fraction of neighbor pairs that share an edge -
making a triangle. For the undirected case, we define this concept formally
below.
###### Definition 7 (Clustering Coefficient)
Let $r$ be the number of edges between nodes with which $v_{i}$ has an edge
and $d_{i}$ be the degree of $v_{i}$. The clustering coefficient,
$C_{i}=\dfrac{2r}{d_{i}(d_{i}-1)}$.
Intuitively, a node with high $C_{i}$ tends to have more pairs of friends that
are also mutual friends. We use the average clustering coefficient as a
network-wide measure of this local property.
Second, we consider modularity ($M$) defined by Newman and Girvan. newman04 .
For a partition of a network, $M$ is a real number in $[-1,1]$ that measures
the density of edges within partitions compared to the density of edges
between partitions. We present a formal definition for an undirected network
below.
###### Definition 8 (Modularity newman04 )
Given partition $C=\\{c_{1},\ldots,c_{q}\\}$, modularity,
$M=\dfrac{1}{2m}\sum_{c\in C}\sum_{i,j\in c}w_{ij}-P_{ij}$
where $m$ is the number of undirected edges; $w_{ij}=1$ if there is an edge
between nodes $i$ and $j$ and $w_{ij}=0$ otherwise;
$P_{ij}=\frac{k_{i}k_{j}}{2m}$
The modularity of an optimal network partition can be used to measure the
quality of its community structure. Though modularity-maximization is NP-hard,
the approximation algorithm of Blondel et al. blondel08 (a.k.a. the “Louvain
algorithm”) has been shown to produce near-optimal partitions.222Louvain
modularity was computed using the implementation available from CRANS at
http://perso.crans.org/aynaud/communities/. We call the modularity associated
with this algorithm the “Louvain modularity.” Unlike the $C$, which describes
local properties, $M$ is descriptive of the community level. For the $31$
networks we considered, $M$ and $C$ appear uncorrelated ($R^{2}=0.0538$,
$p=0.2092$).
We plotted the initial seed set size ($S$) (from our algorithm - averaged over
the $10$ threshold settings) as a function of $M$ and $C$ (Figure 14a) and
uncovered a correlation (planar fit, $R^{2}=0.8666$, $p=5.666\cdot 10^{-13}$,
see Figure 14 A). The majority of networks in Category C (less susceptible to
spreading) were characterized by relatively large $M$ and $C$ (Category C
includes the top nine networks w.r.t. $C$ and top five w.r.t. $M$). Hence,
networks with dense, segregated, and close-knit communities (large $M$ and
$C$) suppress spreading. Likewise, those with low $M$ and $C$ tended to
promote spreading. Also, we note that there were networks that promoted
spreading with dense and segregated communities, yet were less clustered (i.e.
Category A networks Friendster and LiveJournal both have $M\geq 0.65$ and
$C\leq 0.13$). Further, some networks with a moderately large clustering
coefficient were also in Category A (two networks extracted from BlogCatalog
had $C\geq 0.46$) but had a relatively less dense community structure (for
those two networks $M\leq 0.33$).
Figure 14: (A) Louvain modularity ($M$) and average clustering coefficient
($C$) vs. the average seed size ($S$). The planar fit depicted is
$S=43.374\cdot M+33.794\cdot C-24.940$ with $R^{2}=0.8666$, $p=5.666\cdot
10^{-13}$. (B) Same plot at (A) except the averages are over the 12
percentage-based threshold values. The planar fit depicted is $S=18.105\cdot
M+17.257\cdot C-10.388$ with $R^{2}=0.816$, $p=5.117\cdot 10^{-11}$. Table 4:
Regression analysis and network-wide measures for the networks in Categories
A, B, and C. Table 5: Regression analysis and network-wide measures for the
non-symmetric networks.
## 5 Related Work
Tipping models first became popular by the works of Gran78 and Schelling78
where it was presented primarily in a social context. Since then, several
variants have been introduced in the literature including the non-
deterministic version of kleinberg (described later in this section) and a
generalized version of jy05 . In this paper we focused on the deterministic
version. In wattsDodds07 , the authors look at deterministic tipping where
each node is activated upon a percentage of neighbors being activated. Dryer
and Roberts Dreyer09 introduce the MIN-SEED problem, study its complexity,
and describe several of its properties w.r.t. certain special cases of
graphs/networks. The hardness of approximation for this problem is described
in chen09siam . The work of benzwi09 presents an algorithm for target-set
selection whose complexity is determined by the tree-width of the graph -
though it provides no experiments or evidence that the algorithm can scale for
large datasets. The recent work of reichman12 proves a non-trivial upper
bound on the smallest seed set.
Our algorithm is based on the idea of shell-decomposition that currently is
prevalent in physics literature. In this process, which was introduced in
Seidman83 , vertices (and their adjacent edges) are iteratively pruned from
the network until a network “core” is produced. In the most common case, for
some value $k$, nodes whose degree is less than $k$ are pruned (in order of
degree) until no more nodes can be removed. This process was used to model the
Internet in ShaiCarmi07032007 and find key spreaders under the SIR epidemic
model in InfluentialSpreaders_2010 . More recently, a “heterogeneous” version
of decomposition was introduced in baxter11 \- in which each node is pruned
according to a certain parameter - and the process is studied in that work
based on a probability distribution of nodes with certain values for this
parameter.
### 5.1 Notes on Non-Deterministic Tipping
We also note that an alternate version of the model where the thresholds are
assigned randomly has inspired approximation schemes for the corresponding
version of the seed set problem.kleinberg ; leskovec07 ; chen10 Work in this
area focused on finding a seed set of a certain size that maximizes the
expected number of adopters. The main finding by Kempe et al., the classic
work for this model, was to prove that the expected number of adopters was
submodular - which allowed for a greedy approximation scheme. In this
algorithm, at each iteration, the node which allows for the greatest increase
in the expected number of adopters is selected. The approximation guarantee
obtained (less than $0.63$ of optimal) is contingent upon an approximation
guarantee for determining the expected number of adopters - which was later
proved to be $\\#P$-hard. chen10 Recently, some progress has been made toward
finding a guarantee Borgs12 . Further, the simulation operation is often
expensive - causing the overall time complexity to be $O(x\cdot n^{2})$ where
$x$ is the number of runs per simulation and $n$ is the number of nodes
(typically, $x>n$). In order to avoid simulation, various heuristics have been
proposed, but these typically rely on the computation of geodesics - an
$O(n^{3})$ operation - which is also more expensive than our approach.
Additionally, the approximation argument for the non-deterministic case does
not directly apply to the original (deterministic) model presented in this
paper. A simple counter-example shows that sub-modularity does not hold here.
Sub-modularity (diminishing returns) is the property leveraged by Kempe et al.
in their approximation result.
### 5.2 Note on an Upper Bound of the Initial Seed Set
Very recently, we were made aware of research by Daniel Reichman that proves
an upper bound on the minimal size of a seed set for the special case of
undirected networks with homogeneous threshold values. reichman12 The proof
is constructive and yields an algorithm that mirrors our approach (although
Reicshman’s algorithm applies only to that special case). We note that our
work and the work of Reichman were developed independently. We also note that
Reichman performs no experimental evaluation of the algorithm.
Given undirected network $G$ where each node $v_{i}$ has degree $d_{i}$ and
the threshold value for all nodes is $k$, Reichman proves that the size of the
minimal seed set can be bounded by $\sum_{i}\min\\{1,\frac{k}{d_{i}+1}\\}$.
For our integer tests, we compared our results to Reichman’s bound. Our seed
sets were considerably smaller - often by an order of magnitude or more. See
Figure 15 for details.
## 6 Conclusion
As recent empirical work on tipping indicates that it can occur in real social
networks,centola10 ; zhang11 our results are encouraging for viral marketers.
Even if we assume relatively large threshold values, small initial seed sizes
can often be found using our fast algorithm - even for large datasets. For
example, with the FourSquare online social network, under majority threshold
($50\%$ of incoming neighbors previously adopted), a viral marketeer could
expect a $297$-fold return on investment. As results of this type seem to hold
for many online social networks, our algorithm seems to hold promise for those
wishing to “go viral.” An important open question to address in future work is
if a similar decomposition-based approach is viable for finding seed sets
under other diffusion models, such as the independent cascade model kleinberg
and evolutionary graph theory lieberman_evolutionary_2005 as well as
probabilistic variants of the tipping model and diffusion processes on multi-
modal networks snops-iclp . Exploring other models can lead to the development
of decomposition algorithms in domains where social behavior is more dynamic
such as cell-phone networks dyag13 ; otherPho .
###### Acknowledgements.
We would like to thank Gaylen Wong (USMA) for his technical support.
Additionally, we would like to thank (in no particular order) Albert-László
Barabási (NEU), Sameet Sreenivasan (RPI), Boleslaw Szymanski (RPI), Patrick
Roos (UMD), John James (USMA), and Chris Arney (USMA) for their discussions
relating to this work. Finally, we would also like to thank Megan Kearl,
Javier Ivan Parra, and Reza Zafarani of ASU for their help with some of the
datasets. The authors are supported under by the Army Research Office (project
2GDATXR042) and the Office of the Secretary of Defense (project
F1AF262025G001). The opinions in this paper are those of the authors and do
not necessarily reflect the opinions of the funders, the U.S. Military
Academy, or the U.S. Army.
Figure 15: Integer threshold values vs. the seed size divided by Reichman’s
upper bound reichman12 the three categories of networks (categories A-C are
depicted in panels A-C respectively). Note that in nearly every trial, our
algorithm produced an initial seed set significantly smaller than the bound -
in many cases by an order of magnitude or more.
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|
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| 2013-09-11T20:40:11 |
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{
"license": "Public Domain",
"authors": "Paulo Shakarian, Sean Eyre, Damon Paulo",
"submitter": "Paulo Shakarian",
"url": "https://arxiv.org/abs/1309.2963"
}
|
1309.2966
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# Sustainable Software Ecosystems for Open Science:
15 Years of Practice and Experience at Kitware
Marcus D. Hanwell Kitware, Inc., 28 Corporate Drive, Clifton Park, NY 12065
USA.
http://www.kitware.com/ Amitha Perera Kitware, Inc., 28 Corporate Drive,
Clifton Park, NY 12065 USA.
http://www.kitware.com/ Wes Turner Kitware, Inc., 28 Corporate Drive, Clifton
Park, NY 12065 USA.
http://www.kitware.com/ Patrick O’Leary Kitware, Inc., 28 Corporate Drive,
Clifton Park, NY 12065 USA.
http://www.kitware.com/ Katie Osterdahl Kitware, Inc., 28 Corporate Drive,
Clifton Park, NY 12065 USA.
http://www.kitware.com/ Bill Hoffman Kitware, Inc., 28 Corporate Drive,
Clifton Park, NY 12065 USA.
http://www.kitware.com/ Will Schroeder Kitware, Inc., 28 Corporate Drive,
Clifton Park, NY 12065 USA.
http://www.kitware.com/
## Introduction
Mathematics is the core language of science, and for centuries it was
necessary to show the mathematical underpinnings of new research as part of
scientific explorations. This lingua franca provided an essential level of
understandability and precision; providing for unambiguous communication and
rigorous verification of scientific claims beyond the inaccuracies of spoken
languages. However, the last few decades have seen an erosion in this
paradigm. The increased reliance of science on complex computational codes and
large data makes the description of all but the most basic research error
prone, impenetrable, and unverifiable.
This issue is not restricted to just one field of science, but is endemic
throughout the broader scientific community and the consequences of opaque
processes and lack of reproducibility are not trivial. Cases of irreproducible
studies and clinical trials have been making headlines, from Bayer Health Care
stopping nearly two-thirds of its target-validation projects because of
inconsistencies with the initially published claims, to global economic policy
being based on a single fundamentally flawed study by Harvard economists.
These costly mistakes can be remedied much earlier, and before key decisions
are made, simply by returning transparency and precision to the process of
publication and review.
The question we must address is how best to reinstate a common language and
what that language should be. We believe that the only practical choice is to
require that disclosures of scientific research based on complex codes and
data use the very same complex codes and data as the common language of
publication. This means that as new studies and new scientific explorations
are undertaken, the data, methods, and software used by the researchers to
arrive at their conclusions must be made available and accessible to other
researchers and the general populace. If this goal is to be realized then the
standard of software engineering in science must be improved, and sustainable
software ecosystems with meaningful credit must be realized. This is not
simply limited to teaching scientists to write code; if sustainable software
projects are to be established in science then issues such as testing,
licensing, and collaboration must be addressed. Sufficient engineering
discipline is required to realize robust foundations that can be extended
through the use of code review, regression testing and proper citation and
attribution of software used in research.
## Vision
Kitware has been developing open-source solutions and laying the building
blocks for open science for over a decade. We contribute to a wide and varied
range of open-source projects in the scientific domain across many disciplines
including supercomputing, scientific data visualization, large data management
and informatics, computational fluid dynamics, medical imaging, chemistry and
computer vision. Kitware’s specific expertise in this area, including business
practices, collaboration patterns, and growth strategies, are discussed. Based
on this experience, we believe practical frameworks for science revolve around
well-defined and open application programming interfaces (APIs) between key
steps in the analysis workflow from raw data to final figures in published
manuscripts.
Although not directly related to the scientific goals, it is important to
consider licensing in the context of scientific communities. It is not enough
to simply make code and data available for free if we wish to see communities
flourish—shared ownership and access must be provided under permissive
licenses wherever possible. These should encourage reuse, derivatives, and
allow those products to be shared. Permissive licenses such as MIT, BSD, and
Apache 2 for software; and CC0 or CC-BY for data should be strongly encouraged
over copyleft, non-commercial and no derivative clauses that can hamper
further reuse of code or data. When a wealth of data and code is available
under permissive licenses that encourage reuse it is much easier to grow
sustainable software ecosystems—suddenly the method a student developed in
their research can be incorporated into a larger framework even after a
student has moved on to the next challenge. This relies upon the code being
published and made available under permissive licenses allowing for reuse in
other contexts. The same path can be considered for data, where individual
findings can be combined into larger collections and analyses performed on
aggregate data to make new discoveries outside the scope of the original
research in many cases.
## Practice and Experience
Kitware started as an experiment in open source education beginning when three
General Electric research employees wanted to author a book on visualization
using a new programming language called C++. The book included code, and in
order to encourage reuse and improve the general standard in the field they
were granted authorization by GE to retain copyright and publish the code with
the book. This would later become the Visualization Toolkit (VTK). In 1998 two
of the authors left GE and with several others founded a small company in
upstate New York named Kitware. Fifteen years later the adventure continues,
with over 100 employees and sustained year-on-year growth.
### Baseline Implementations
As noted, VTK began life as code to support a book. The goal was to improve
the state-of-the-art in the visualization community by putting tested and
verified implementations of important visualization and analysis algorithms
into a software library that was permissively licensed and available for use
by a broad and open community. The intent was to provide these implementations
as a concrete instantiation of complex algorithms so that researchers around
the world did not have to do this from the papers of the time, with varying
degrees of success and different levels of verification. This same goal lies
at the heart of the Insight Toolkit in medical imaging for image segmentation
and registration. Over the years things such as the Insight Journal were added
so that new algorithms could be proposed in formal publications that included
code, data, and baselines demonstrating the algorithm. If accepted, the
algorithm would be merged into the main codebase and made available to the
wider community.
### Business Model
The majority of well-known software companies use a licensing and intellectual
property model in order to derive revenue from their development activity. In
the sciences this is often bolstered by funding from major agencies as well.
Kitware generates no intellectual property, and has moved away from the
licensing model to a large extent—focusing our business activity on a services
model. This is an extremely effective strategy when working in high-
performance computing, and addresses one of the major flaws in many open-
source projects where support is difficult to find. Services make up a larger
share of the industry sector, with companies such as IBM and Red Hat deriving
most of their income from services rather than licensing. This model enables
us to work with academics, national laboratories, and industry as partners
often with joint funding streams to develop new features and make headway in
different avenues of research and development. Through the use of permissive
licensing models, rigorous software processes and agile development
methodologies we have been able to make significant progress and received
recognition for major projects such as VTK and ParaView.
### Developing Communities
One of our core focus areas at Kitware is software process and developing
robust communities around software projects that can grow beyond individual
contracts. This process is shown at a high level in Figure 1 where developers
commit code, it is automatically tested, dashboards gather testing information
and disseminate that information. This process is powered by CMake, which
began as a build system for ITK when existing build systems proved inadequate,
with the simple goal of building C++ code on all major platforms. This was
later augmented with CTest that would run and report the results of automated
testing to Dart, and later CDash. Packaging was addressed using CPack,
enabling some degree of abstraction when creating binary installers which are
now routinely generated each night so that non-developers can install the
latest build of a software project to see if a bug has been fixed.
Figure 1: Overview of the Kitware software process.
Kitware has built up software workflows around version control systems such as
RCS, CVS, Subversion and more recently Git (among other distributed version
control systems). These processes are designed to minimize the barrier to
entry while maintaining the high quality of the software using automated
repository hooks, continuous integration, and advanced testing capabilities
such as image-based regression images to verify that visualization algorithms
produce consistent and correct output on a large number of platforms beyond
the capability of any individual contributor.
Any community needs to establish a strong process to remain viable over the
long term. This requires a mixture of technical and social resources to
facilitate productive engagement and growth. Scientific projects often have
limited resources, with a high turnover of developers/engineers as they move
through their academic careers. This means that it is essential to establish
procedures to help new community members get up to speed, and to retain group
knowledge as people move on to new positions. Some will maintain long-term
ties through different positions, as is normal in more mainstream open-source
projects, others will not and it is important to ensure the parts of the
project they developed do not atrophy.
### Open Software Process
Kitware engineers and researchers have experimented with a number of
techniques to improve induction of new community members and employees into
our company and software communities, to support open-source development, and
the growth of our projects. Most importantly, these experiments are more than
just an academic exercise. All of the techniques are applied as part of our
ongoing, commercial, software business. The techniques that offer pragmatic
benefits to our business and open-source communities survive. Figure 2 shows a
more detailed view of the software process used by the Visualization Toolkit
(VTK).
Figure 2: VTK’s software process incorporating online code review, automated
testing of topic branches, and nightly testing with distributed mirrors.
Recently we started using Gerrit to conduct online code reviews, and added
automated triggers so that proposed changes are built on all major operating
systems. This simple modification allows us to highlight potential platform-
specific issues before engineers review code, saving scarce engineering time.
Our software projects are developed using test-driven development
methodologies, and require the submission of tests for new code or features
ensuring that these features are automatically tested and verified before they
are merged.
We have a strong commitment to testing. We run nightly dashboards to provide
wider coverage, ensuring that any problems are normally highlighted within 24
hours of merging them into the repository and allowing developers to review
and remedy regressions while the changes are still fresh in their mind. All of
these processes are learning opportunities, enabling developers to quickly
discover what will and won’t work with different environments as new features
are developed. These advantages were only made possible in recent years by
migrating to distributed version control systems that enable development
outside of the main repository and for changes to be pushed to staging
locations before being merged into the main development branches.
We have worked hard to close the gap between developers of the frameworks and
applications and their users by automatically generating binary installers.
This enables them to download the latest nightly binary of any given project
to check if their bug still exists, or work more closely with developers on
exercising new features before the final release is made. Underlying all of
this are multiple communication mediums including mailing lists, wikis, and
bug trackers to offer a comprehensive set of community tools. Having a
commercial company backing up these projects also offers the ability to
provide books, webinars, and on-site courses.
### Mature Communities
Kitware hosts or plays a leading role in the maintenance and engineering on a
number of projects, and over the years has learned to be adaptable. The oldest
Kitware-led community is VTK, which also has one of the longest running
version histories of any project whose history is publicly available. This is
a community that grew up organically around the VTK book, and many pieces of
the code were experimented with and added as the community grew. According to
Ohloh the project has had nearly 55,000 commits made by over 200 contributors
and includes more than 3.5 million lines of code. The first recorded VTK
commit occurred nearly 20 years ago and development remains active. Even when
looking at the last 12 months the project has seen over 3,000 commits from 66
contributors under a permissive BSD license. Figure 3 shows a graphical
overview of the project’s activity from 1994 to present.
Figure 3: Statistics generated by Ohloh showin lines of code, and commit rates
for VTK.
Even so, these numbers understate the total significance. For instance, they
do not count the number of commercial contracts where Kitware supported the
toolkit by providing engineers to extend VTK, or create/augment focused
applications that leverage VTK, or research conducted using VTK. The
development is also funded by several ongoing contracts with established
groups, and funding streams from various federal agencies such as the DOE,
NIH, DOD, DARPA, and NSF. Due to VTK’s importance to researchers in the
medical area, it was granted a rare R01 maintenance grant from the NIH. Under
this contract a consortium of companies and universities led by Kitware is
working to update and overhaul aspects of the VTK rendering subsystems and to
make other improvements over the next four years. VTK was able to garner a
large number of support letters from the community of users and developers
making use of VTK in their research, applications and products.
After the success of VTK, the ITK project was, at least in part, founded to
emulate the success of VTK in the medical segmentation and registration field.
ITK receives ongoing funding from the National Library of Medicine, initially
to provide algorithms that could operate on and make sense of the data from
the visible human project. It is a little younger, with its first commit in
2000, and is interesting as it was one of the first major projects that began
after Kitware’s founding and was designed from the ground up as an open-source
project. It has seen nearly 45,000 commits in that time from over 200
contributors with over 1.5 million lines of code. Over the last year more than
1300 commits were made by 67 contributors under a permissive Apache 2 license.
The ITK project also has an associated non-profit organization that is used to
coordinate community efforts and manage some of the funding and development.
The history of some of these older projects that have seen one or two decades
of development effort should be contrasted with younger projects at different
stages of development. Experimentation is encouraged, and not all of the
projects that are started will achieve the critical mass required to become
successful open source projects, but our established processes seek to make
experimentation easy and the permissive open-source licensing means that even
“failed” projects will often live on in other forms if the code developed for
the project is useful.
A brief summary of some of these projects include Arbor/Avatol—an NSF funded
project where Kitware is an engineering subcontractor in a larger project led
by the University of Idaho. The goals of this project are to provide tools for
the analysis and exploration of the Tree of Life. The XDATA project is a new
project funded by DARPA with a large consortium (led by Kitware as the prime
contractor) tasked with creating open-source technology to address big data
analytics. The Open Chemistry project was funded by a DOD SBIR awarded to
Kitware, with the aim of creating a suite of desktop applications for
computational chemists.
In addition to creating, growing and shepherding our own communities Kitware
has more recently begun helping other organizations to build successful
communities. This includes the VA with the OSEHRA project, NA-MIC, and caBIG,
all of whom use Kitware’s expertise to assist them in growing vibrant and
successful communities around their projects.
### Concluding Remarks
Sustainable software ecosystems are difficult to build, and require concerted
effort, community norms and collaborations. In science it is especially
important to establish communities in which faculty, staff, students and open-
source professionals work together and treat software as a first-class product
of scientific investigation—just as mathematics is treated in the physical
sciences. Kitware has a rich history of establishing collaborative projects in
the science, engineering and medical research fields, and continues to work on
improving that model as new technologies and approaches become available. This
approach closely follows and is enhanced by the movement towards practicing
open, reproducible research in the sciences where data, source code,
methodology and approach are all available so that complex experiments can be
independently reproduced and verified.
### License
This document is released under the Creative Commons Attribution 3.0 license
(CC-BY), see http://creativecommons.org/licenses/by/3.0/.
|
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| 2013-09-11T20:52:56 |
2024-09-04T02:49:50.769299
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Marcus D. Hanwell, Amitha Perera, Wes Turner, Patrick O'Leary, Katie\n Osterdahl, Bill Hoffman and Will Schroeder",
"submitter": "Marcus Hanwell",
"url": "https://arxiv.org/abs/1309.2966"
}
|
1309.2975
|
These are not the k-mers you are looking for: efficient online k-mer counting
using a probabilistic data structure
Qingpeng Zhang1, Jason Pell1, Rosangela Canino-Koning1, Adina Chuang Howe2,3,
C. Titus Brown1,2∗
1 Department of Computer Science and Engineering, Michigan State University,
East Lansing, MI, USA
2 Department of Microbiology and Molecular Genetics, Michigan State
University, East Lansing, MI, USA
3 Department of Plant, Soil, and Microbial Sciences, Michigan State
University, East Lansing, MI, USA
$\ast$ E-mail: [email protected]
## Abstract
K-mer abundance analysis is widely used for many purposes in nucleotide
sequence analysis, including data preprocessing for de novo assembly, repeat
detection, and sequencing coverage estimation. We present the khmer software
package for fast and memory efficient online counting of k-mers in sequencing
data sets. Unlike previous methods based on data structures such as hash
tables, suffix arrays, and trie structures, khmer relies entirely on a simple
probabilistic data structure, a Count-Min Sketch. The Count-Min Sketch permits
online updating and retrieval of k-mer counts in memory which is necessary to
support online k-mer analysis algorithms. On sparse data sets this data
structure is considerably more memory efficient than any exact data structure.
In exchange, the use of a Count-Min Sketch introduces a systematic overcount
for k-mers; moreover, only the counts, and not the k-mers, are stored. Here we
analyze the speed, the memory usage, and the miscount rate of khmer for
generating k-mer frequency distributions and retrieving k-mer counts for
individual k-mers. We also compare the performance of khmer to several other
k-mer counting packages, including Tallymer, Jellyfish, BFCounter, DSK, KMC,
Turtle and KAnalyze. Finally, we examine the effectiveness of profiling
sequencing error, k-mer abundance trimming, and digital normalization of reads
in the context of high khmer false positive rates. khmer is implemented in C++
wrapped in a Python interface, offers a tested and robust API, and is freely
available under the BSD license at github.com/ged-lab/khmer.
## Introduction
The goal of k-mer counting is to determine the number of occurrences for each
fixed-length word of length k in a DNA data set [1]. Efficient k-mer counting
plays an important role in many bioinformatics approaches, including data
preprocessing for de novo assembly, repeat detection, and sequencing coverage
estimation [2].
Short-read shotgun sequencing data is both relatively sparse in k-mers and
contains many erroneous k-mers. For typical values of k such as 32 these data
sets are sparse, as only a small fraction of the total possible number of
k-mers ($4^{32}$) are actually present in any genome or read data sets derived
from the genome. The high error rate (e.g. Illumina has a 0.1-1% per-base
error rate [3]) generates many unique k-mers. As the total number of generated
reads increases, the total number of errors grows with it linearly. This leads
to data sets where the erroneous k-mers vastly outnumber the true k-mers [4].
Tracking and counting the resulting large number of k-mers, most of which are
erroneous, has become an unavoidable and challenging task in sequence analysis
[5].
A variety of k-mer counting approaches, and standalone software packages
implementing them, have emerged in recent years; this includes Tallymer,
Jellyfish, BFCounter, DSK, KMC, Turtle and KAnalyze [2, 1, 6, 7, 8, 9, 10].
These approaches and implementations each offer different algorithmic trade-
offs and enable a non-overlapping set of functionality. Tallymer uses a suffix
tree to store k-mer counts in memory and on disk [2]. Jellyfish stores k-mer
counts in in-memory hash tables, and makes use of disk storage to scale to
larger data sets [1]. BFCounter uses a Bloom filter as a pre-filter to avoid
counting unique k-mers, and is the first published probabilistic approach to
k-mer counting [6]. DSK adopts an approach to k-mer counting that enables
time- and memory-efficient k-mer counting with an explicit trade-off between
disk and memory usage [7]. KMC and KAnalyze rely primarily on fast and
inexpensive disk access to count k-mers in low memory [8, 10]. Turtle provides
several different containers that offer different false positive and false
negative tradeoffs when counting k-mers [9].
Our motivation for exploring efficient k-mer counting comes from our work with
metagenomic data, where we routinely encounter data sets that contain
$300\times 10^{9}$ bases of DNA and over 50 billion distinct k-mers [11]. To
efficiently filter, partition, and assemble these data, we need to store
counts for each of these k-mers in main memory, and query and update them in
realtime — a set of functionality not readily offered by current packages.
Moreover, we wish to enable the use of cloud and desktop computers, which may
have poor I/O performance or limited memory. These needs have dictated our
exploration of efficient in-memory k-mer counting techniques.
Below, we describe an implementation of a simple probabilistic data structure
for k-mer counting. This data structure is based on a Count-Min Sketch [12], a
generalized probabilistic data structure for storing the frequency
distributions of distinct elements. Our implementation extends an earlier
implementation of a Bloom filter [13], which has been previously used in
bioinformatics applications, such as sequence matching [14], k-mer counting
[6], and de Bruijn graph storage and traversal [15, 16]. Many other variations
of Bloom filters have been proposed [17], including counting Bloom filters
[18], multistage filters [19], and spectral Bloom filters [20], which are
related to the Count-Min Sketch and our khmer implementation.
Probabilistic approaches can be particularly memory efficient for certain
problems, with memory usage significantly lower than any exact data structure
[15]. However, their use introduces set membership or counting false
positives, which have effects that must be analyzed in the context of specific
problems. Moreover, unlike existing techniques, the Count-Min Sketch stores
only counts; k-mers must be retrieved from the original data set. In exchange,
the low memory footprint enabled by this probabilistic approach enables online
updating and retrieval of k-mer counts entirely in memory, which in turn
supports streaming applications such as digital normalization [21].
We use the Amazon cloud to compare time, memory, and disk usage of our k-mer
counting implementation with that of other k-mer counting software packages,
for two problems. First, we generate a k-mer abundance distribution for large
data sets; and second, we query many individual k-mer counts at random from a
previously constructed k-mer count database. We show that khmer is competitive
in speed, memory, and disk usage for these problems. We also analyze the
effects of counting error on calculations of the k-mer count in sequencing
data sets, and in particular on metagenomic data sets. Finally, we discuss
khmer’s miscount performance in the context of two specific applications: low-
abundance k-mer trimming of reads, and digital normalization.
The khmer software [22] is implemented in C++ in a Python wrapper, enabling
flexible use and reuse by users with a wide range of computational expertise.
The software package is freely available for academic and commercial use and
redistribution under the BSD license at github.com/ged-lab/khmer/. khmer comes
with substantial documentation and many tutorials, and contains extensive unit
tests. Moreover, we have built several applications on top of khmer, including
memory-efficient de Bruijn graph partitioning [15] and lossy compression of
short-read data sets for assembly [21].
## Results
### Implementing a Count-Min Sketch for k-mers
The two basic operations supported by khmer are increment_count(kmer) and c =
get_count(kmer). Both operate on the data structure in memory, such that
neither incrementing a count nor retrieving a count involves disk access.
The implementation details are similar to those of the Bloom filter in [15],
but with the use of 8 bit counters instead of 1 bit counters. Briefly, Z hash
tables are allocated, each with a different size of approximately H bytes
($H_{1},H_{2},...,H_{Z}$); the sum of these hash table sizes must fit within
available main memory. To increment the count for a particular k-mer, a single
hash is computed for the k-mer, and the modulus of that hash with each hash
table’s size H gives the location for each hash table; the associated count in
each hash table is then incremented by 1. We use different sizes for each hash
table so as to vary the hash function. Even if two k-mers have the same
modulus in one hash table (a collision), they are unlikely to collide in the
other hash tables. To retrieve the count for a k-mer, the same hash is
computed and the minimum count across all hash tables is computed. While
different in implementation detail from the standard Bloom filter, which uses
a single hash table with many hash functions, the performance details are
identical [15]. One particularly important feature of the Count-Min Sketch is
that the counting error is one-sided [12]. Because counts are only
incremented, collisions result in inflated miscounts; if there is no collision
for a particular k-mer, the count is correct.
An additional benefit of the Count-Min Sketch is that it is extremely easy to
implement correctly, needing only about 3 dozen lines of C++ code for a simple
threadsafe implementation. (We have described how khmer scales with multiple
threads in [23].)
To determine the expected false positive rate — the average frequency with
which a given k-mer count will be incorrect when retrieved — we can look at
the hash table load. Suppose N distinct k-mers have been counted using Z hash
tables, each with size H. The probability that no collisions happened in a
specific entry in one hash table is $(1-1/H)^{N}$, or approximately
$e^{-N/H}$. The individual collision rate in one hash table is then $\approx
1-e^{-N/H}$. The total collision rate, which is the probability that a
collision occurred in each entry where a k-mer maps across all Z hash tables,
is $\approx(1-e^{-N/H})^{Z}$, which is also the expected false positive rate.
While the false positive rate can easily be calculated from the hash table
load, the average miscount — the degree to which the measured count differs
from the true count — depends on the k-mer frequency distribution, which must
be determined empirically. We analyze the effects of this below.
### Choosing number and size of hash tables used for k-mer counting
The false positive rate depends on the number of distinct k-mers $N$, the
number of hash tables $Z$, and the size of the hash tables $H$:
$f\approx(1-e^{-N/H})^{Z}$, with an associated memory usage of $M=HZ$. We face
two common scenarios: one in which we have a fixed number of k-mers $N$ and
fixed memory $M$ and we want to calculate the optimal number of hash tables
$Z$; and one in which we have a desired maximum false positive rate $f$ and a
fixed number of k-mers $N$, and we want to calculate the minimum memory usage
required to achieve $f$.
For fixed memory $M$ and number of distinct k-mers $N$, the optimal number of
hash tables can be found by minimizing $f$; taking the derivative, $df/dZ$,
with $f\approx\exp(Z\log(1-e^{-ZN/M}))$ and solving for 0, we find that $f$ is
minimized when $Z=\log(2)*(M/N)$ (see [24] for details).
Given a desired false positive rate $f$ and a fixed number of k-mers $N$, the
optimal memory usage can be calculated as follows. First, the optimal number
of hash tables is determined by the expected false positive rate alone:
$Z=\log_{0.5}f$. Using this $Z$, the minimum average hash table size $H$
necessary to achieve $f$ can be calculated as $H=(\log_{0.6185}(f)\times N)/Z$
(see [24] for details).
A remaining problem is that the number of distinct k-mers $N$ is typically not
known. However, memory- and time-efficient algorithms for calculating $N$ do
exist and we plan to implement this in khmer in the future [25].
### khmer efficiently calculates k-mer abundance histograms
We measured time and memory required to calculate k-mer abundance histograms
in five soil metagenomic read data sets using khmer, Tallymer, Jellyfish, DSK,
KMC, Turtle, and KAnalyze (Table 1; Figures 1 and 2). We chose to benchmark
abundance histograms because this functionality is common to all the software
packages, and is a common analysis approach for determining assembly
parameters [26]. We applied each package to increasingly large subsets of a
50m read soil metagenome data set [11]. For the BFCounter, KMC, Turtle and
KAnalyze packages, which do not generate k-mer abundance distribution
directly, we output the frequency of each k-mer to a file but do no further
analysis.
Figure 1 shows that the time usage of the khmer approach is comparable to DSK
and BFCounter, and, as expected, increases linearly with data set size.
Tallymer is the slowest of the four tools in this testing, while KMC, Turtle,
and Jellyfish are the fastest.
From Figure 2, we see that the memory usage of Jellyfish, Tallymer, BFCounter,
and Turtle increases linearly with data set size. Tallymer uses more memory
than Jellyfish generally, while BFCounter and Turtle have considerably lower
memory usage. DSK, KMC, and KAnalyze use constant memory across the data sets,
but at the cost of more limited functionality (discussed below).
The memory usage of khmer also increases linearly with data set size as long
as we hold the false positive rate constant. However, the memory usage of
khmer varies substantially with the desired false positive rate: we can
decrease the memory usage by increasing the false positive rate as shown in
Figure 2. We also see that with a low false positive of 1%, the memory usage
is competitive with Tallymer and Jellyfish; with a higher 5% false positive
rate, the memory usage is lower than all but the disk-based DSK; with an false
positive rate as high as 20%, the memory usage is further lower, close to DSK,
KAnalyze, and KMC.
We also measured disk usage during counting. Figure 3 shows that the disk
usage also increases linearly with the number of k-mers in the data set. For a
high-diversity metagenomic data set of 5 GB, the disk usage of both Jellyfish
and Tallymer is around 30 GB. khmer counts k-mers entirely in working memory
and does not rely on any on-disk storage to store or retrieve k-mer counts,
although for practicality the hash tables can be saved for later reuse; the
uncompressed disk usage for khmer in Figure 3 is the same as its memory. At
the expense of more time, khmer supports saving and loading gzip-compressed
hash tables, which are competitive in size to DSK’s on-disk database (Figure
3, dashed line).
### khmer accesses k-mer counts efficiently
We measured the time it took to access 9.7m 22-mers across five different data
sets after the initial databases had been built (Figure 4). Note that
Tallymer, Jellyfish, and khmer all support random access to k-mer counts,
while BFCounter, DSK, KMC, Turtle and KAnalyze do not. Here, khmer performed
well, dramatically outperforming Jellyfish and Tallymer. In all three cases,
system time dominated the overall time required to retrieve k-mers, suggesting
that the primary reason for the increase in retrieval time was due to the
increased size of the database on the disk (data not shown). In particular,
khmer is independent of the size of the database in retrieval time once the
hash tables are loaded into memory.
### The measured counting error is low on short-read data
Due to the use of Count-Min Sketch and its lack of collision tracking, khmer
will report some incorrect counts for k-mers; these counts are always higher
than the true counts, up to the bound of 255 (a limit imposed by our use of
8-bit counters). The frequency with which incorrect counts are reported can be
estimated from the hash table load. However, the expected miscount — the
difference between the true k-mer frequency and the reported k-mer frequency —
cannot be calculated without knowing the distribution of k-mer abundances; in
general, the average miscount will be small if the data is left-skewed. As
noted by Melsted and Pritchard, a large number of k-mers in short-read data
are low-abundance, leading to precisely the skew that would yield low
miscounts [6]. Here we use both real and simulated data sets to evaluate the
counting performance in practice.
Figure 5 shows the relationship between average miscount and counting false
positive rate for five different test data sets with similar numbers of
distinct k-mers: one metagenome data set; a simulated set of random k-mers; a
simulated set of reads, chosen with 3x coverage and 1% error; a simulated set
of reads (3x) with no error; and a set of E. coli reads (Table 2). Even when
the counting false positive rate is as high as 0.9 — where 90% of k-mers have
an incorrect count — the average miscount is still below 4.
We separately analyzed the average percentage miscount between true and false
k-mers; e.g. an miscount of 4 for a k-mer whose true count is 1 would be 400%.
Figure 6 shows the relationship between average miscount and counting false
positive rate for the same five data sets as in Figure 5. For a false positive
rate of 0.1 (10% of k-mer counts are incorrect), the average percentage
miscount is less than 10% for all five data sets; this will of course
generally be true, because the average miscount is bounded by the product of
the false positive rate with k-mer abundance.
We see here that for a fixed false positive rate, the simulated reads without
error have the highest average miscount, and the randomly generated k-mers
have the lowest average miscount. This is because these two abundance
distributions have the least and most left-skew, respectively: the simulated
reads without error have no abundance-1 k-mers, while the randomly generated
k-mers are entirely low abundance.
### Sequencing error profiles can be measured with k-mer abundance profiles
One specific use for khmer is detecting random sequencing errors by looking at
the k-mer abundance distribution within reads [27]. This approach, known also
as “k-mer spectral analysis”, was first proposed in by [28] and further
developed in [29]. The essential idea is that low-abundance k-mers contained
in a high-coverage data set typically represent random sequencing errors.
A variety of read trimming and error correcting tools use k-mer counting to
reduce the error content of the read data set, independent of quality scores
or reference genomes [30]. This is an application where the counting error of
the Count-Min Sketch approach used by khmer may be particularly tolerable: it
will never falsely call a high-abundance k-mer as low-abundance because khmer
never underestimates counts.
In Figure 7, we use khmer to examine the sequencing error pattern of a 5m-read
subset of an Illumina reads data set from single-colony sequencing of E. coli
[31]. The high rate of occurrence of unique k-mers close to the 3’ end of
reads is due to the increased sequencing error rate at the 3’ end of reads.
### khmer can be applied iteratively to read trimming
We next evaluated the effect of false-positive induced miscounts on read
trimming, in which reads are truncated at the first low-abundance k-mer.
Because the Count-Min Sketch never undercounts k-mers, reads will never be
erroneously trimmed at truly high-abundance k-mers; however, reads may not be
trimmed correctly when miscounts inflate the count of low-abundance k-mers. In
cases where many errors remain, read trimming can potentially be applied
multiple times, with each round reducing the total number of k-mers and hence
resulting in lower false positive rates for the same memory usage.
We performed six iterations of unique k-mer trimming on 5 million Illumina
reads from sequencing of E. coli, with memory usage less than 30 MB. For each
iteration we measured empirical false positive rate compared with number of
bases trimmed as well as the total number of k-mers (Table 3). In the first
round, the estimated false positive rate was 80.0%, and 13.5% of the total
bases were removed by trimming reads at low-abundance k-mers; the second
iteration had a false positive rate of 37.7%, and removed only 1.5% additional
data; and by the fourth iteration the false positive rate was down to 23.2%
with 0.0% of the data removed.
The elimination of so many unique k-mers (column 5) in the first pass was
unexpected: the high false positive rate should have resulted in fewer k-mers
being identified as unique, were the erroneous k-mers independent of each
other. Upon examination, we realized that in Illumina data erroneous k-mers
typically come from substitution errors that yield runs of up to k erroneous
k-mers in a row [30]. When trimming reads with high false positive rates,
these runs are typically trimmed after the first few unique k-mers, leaving
unique k-mers at the 3’ end. Because of this we hypothesized that high-FP rate
trimming would result in the retention of many unique k-mers at the 3’ end of
the read, and this was confirmed upon measurement (Table 3, column 6, pass 1
vs pass 2).
In comparison to quality-based trimming software such as seqtk and FASTX,
trimming at unique k-mers performed very well: in this data set, all unique
k-mers represent errors, and even with an initial false positive rate of 80%,
khmer outperformed all but the most stringent seqtk run (Table 3). With a
lower false positive rate or multiple passes, khmer eliminates more erroneous
k-mers than seqtk or FASTX. The tradeoff here is in memory usage: for larger
data sets, seqtk and FASTX will consume the same amount of memory as on
smaller data sets, while khmer’s memory usage will need to grow with the data
set size.
### Using khmer for digital normalization, a streaming algorithm
Digital normalization is a lossy compression algorithm that discards short
reads based on saturating coverage of a de Bruijn graph [21]. While several
non-streaming implementations exist, including Trinity’s in silico
normalization [32, 33], digital normalization can be efficiently implemented
as a streaming algorithm. In the streaming implementation, if a read is not
kept, it is not loaded into the Count-Min Sketch structure, and the false
positive rate does not increase. For high coverage data sets, the digital
normalization algorithm is sublinear in memory because it does not collect the
majority of k-mers in those data sets [21]. This has the advantage of enabling
low-memory preprocessing of both high-coverage genomic data sets, as well as
mRNAseq or metagenomic data sets with high-coverage components [21, 11].
While digital normalization is already implemented inside khmer, previous work
did not explore the lower bound on memory usage for effective digital
normalization. In particular, the effects of high false positive rates have
not been examined in any prior work.
We applied digital normalization to the E. coli data set used above, and chose
seven different Count-Min Sketch sizes to yield seven different false positive
rates 4. The data set was normalized to a k-mer coverage of 20 and the
resulting data were evaluated for retention of true and erroneous k-mers, as
in [21] (Table 4). The results show that digital normalization retains the
same set of underlying “true” k-mers until the highest false positive rate of
100% (Table 4, column 5), while discarding only about 2% additional reads
(Table 4, column 6).
To evaluate the effect of digital normalization with high false positive rates
on actual genome assembly, we next performed normalization to a coverage of 20
with the same range of false positive rates as above. We then assembled this
data with Velvet [34] and compared the resulting assemblies to the known E.
coli MG1655 genome using QUAST (Table 5). To our surprise, we found that even
after executing digital normalization with a false positive rate of 83.2%, a
nearly complete assembly was generated. No progressive increase in
misassemblies (measured against the real genome with QUAST) was seen across
the different false positive rates (data not shown). This suggests that below
83.2% FP rate, the false positive rate of digital normalization has little to
no effect on assembly quality with Velvet. (Note that the Velvet assembler
itself used considerably more memory than digital normalization.)
While these results are specific to Velvet and the coverage parameters used in
digital normalization, they do suggest that no significant information loss
occurs due to false positive rates below 80%. Further evaluation of assembly
quality in response to different normalization parameters and assemblers is
beyond the scope of of this paper.
## Discussion
### khmer enables fast, memory-efficient online counting
khmer enables memory- and time-efficient online counting (Figures 1, 2, and
4). This is particularly important for the streaming approaches to data
analysis needed as data set sizes increase. Because query and updating of
k-mer counts can be done directly as data is being loaded, with no need for
disk access or an indexing step, khmer can also perform well in situations
with poor disk I/O performance. (Note that BFCounter also supports online
k-mer counting [6].)
### khmer is a generally useful k-mer counting approach
In addition to online counting, khmer offers a general range of useful
performance tradeoffs for disk I/O, time and memory. From the performance
comparison between khmer and other k-mer counting packages in calculating
k-mer abundance distributions, khmer is comparable with existing packages. In
time, khmer performs competitively with DSK and BFCounter (Figure 1); khmer
also provides a way to systematically trade memory for miscounts across a wide
range of parameters (Figure 2). khmer’s uncompressed disk storage is
competitive with Jellyfish, and, in situations where disk space is at a
premium, khmer can take advantage of gzip compression to provide storage
similar to that of DSK (Figure 3, purple line with boxes).
KMC, DSK, and KAnalyze perform especially well in memory usage for calculating
the abundance distribution of k-mers. However, in exchange for this
efficiency, retrieving specific k-mer counts at random is likely to be quite
slow, as DSK is optimized for iterating across partition sets of k-mers rather
than randomly accessing k-mer counts.
For retrieving the counts of individual k-mers, khmer is significantly faster
than both Tallymer and Jellyfish. This is not surprising, since this was a
primary motivation for the development of khmer.
### khmer memory usage is fixed and low
The memory usage of the basic Count-Min Sketch approach is fixed: khmer’s
memory usage does not increase as data is loaded. While this means that khmer
will never crash due to memory limitations, and all operations can be
performed in main memory without recourse to disk storage, the false positive
rate may grow too high. Therefore the memory size must be chosen in light of
the false positive rate and miscount acceptable for a given application. In
practice, we recommend choosing the maximum available memory, because the
false positive rate decreases with increasing memory and there are no negative
effects to minimizing the false positive rate.
For any given data set, the size and number of hash tables will determine the
accuracy of k-mer counting with khmer. Thus, the user can control the memory
usage based on the desired level of accuracy (Figure 2). The time usage for
the first step of k-mer counting, consuming the reads, depends on the total
amount of data, since we must traverse every k-mer in every read. The second
step, k-mer retrieval, is algorithmically constant for fixed k; however, for
practicality, the hash tables are usually saved to and loaded from disk,
meaning that k-mer retrieval time depends directly on the size of the database
being queried.
The memory usage of khmer is particularly low for sparse data sets, especially
since only main memory is used and no disk space is necessary beyond that
required for the read data sets. This is no surprise: the information
theoretic comparison in [15] shows that, for sparse sequencing data sets,
Bloom filters require considerably less memory than any possible exact
information storage for a wide range of false positive rates and data set
sparseness.
In our implementation we use 1 byte to store the count of each k-mer in the
data structure. Thus the maximum count for a k-mer will be 255. In cases where
tracking bigger counts is required, khmer also provides an option to use an
STL map data structure to store counts above 255, with the trade-off of
significantly higher memory usage. In the future, we may extend khmer to
counters of arbitrary bit sizes.
### False positive rates in k-mer counting are low and predictable
The Count-Min Sketch is a probabilistic data structure with a one-sided error
that results in random overestimates of k-mer frequency, but does not generate
underestimates.
In the Count-Min Sketch, the total memory usage is fixed; the memory usage,
the hash functions, and the total number of distinct objects counted all
influence the accuracy of the count. While the probability of an inaccurate
count can easily be estimated based on the hash table load, the miscount size
is dependent on details of the frequency distribution of k-mers [12].
More specifically, in the analysis of the Count-Min Sketch, the difference
between the incorrect count and actual count is related to the total number of
k-mers in a data set and the size of each hash table [12]. Further study has
shown that the behavior of Count-Min Sketch depends on specific
characteristics of the data set under consideration, especially left-skewness
[35, 36]. These probabilistic properties suit short reads from next generation
sequencing data sets: the miscounts are low because of the highly left-skewed
abundance distribution of k-mers in these data sets.
Figures 5 and 6 demonstrate these properties well. We see more correct
counting for error-prone reads from a genome than for error-free reads from a
genome, with a normal distribution of k-mer abundance. Thus, this counting
approach is especially suitable for high diversity data sets, such as
metagenomic data, in which a larger proportion of k-mers are low abundance or
unique due to sequencing errors.
### Real-world applications for khmer
For many applications, an approximate k-mer count is sufficient. For example,
when eliminating reads with low abundance k-mers, we can tolerate a certain
number of low-frequency k-mers remaining in the resulting data set falsely. If
RAM-limited we can do the filtering iteratively so that at each step we are
making more effective use of the available memory.
In practice, we have found that a false positive rate of between 1% and 10%
offers acceptable miscount performance for a wide range of tasks, including
error profiling, digital normalization and low-abundance read-trimming.
Somewhat surprisingly, false positive rates of up to 80% can still be used for
both read trimming and digital normalization in memory-limited circumstances,
although multiple passes across the data may be needed.
For many applications, the fact that khmer does not break an imposed memory
bound is extremely useful, since for many data sets — especially metagenomic
data sets — high memory demands constrain analysis [11, 37]. Moreover, because
the false positive rate is straightforward to measure, the user can be warned
that the results should be invalidated when too little memory is used. When
combined with the graceful degradation of performance for both error trimming
and digital normalization, khmer readily enables analysis of extremely large
and diverse data sets [38]. In an experiment to assemble the reads of a soil
metagenomic sample collected from Iowa prairie, the number of reads to
assemble drops from 3.3 million to 2.2 million and the size of the data set
drops from 245GB to 145GB accordingly after digital normalization [11]. 240GB
memory was used in the process. This also shows that khmer works well to
analyze large, real-world metagenomic data sets.
### Conclusion
K-mer counting is widely used in bioinformatics, and as sequencing data set
sizes increase, graceful degradation of data structures in the face of large
amounts of data has become important. This is especially true when the
theoretical and practical effects of the degradation can be predicted (see
e.g. [6, 15, 9]). This is a key property of the Count-Min Sketch approach, and
its implementation in khmer.
The khmer software implementation offers good performance, a robust and well-
tested Python API, and a number of useful and well-documented scripts. While
Jellyfish, DSK, KMC, and Turtle also offer good performance, khmer is
competitive, and, because it provides a Python API for online counting, is
flexible. In memory-limited situations with poor I/O performance, khmer is
particularly useful, because it will not break an imposed memory bound and
does not require disk access to store or retrieve k-mer counts. However, in
exchange for this memory guarantee, counting becomes increasingly incorrect as
less memory is used or as the data set size grows large; in many situations
this may be an acceptable tradeoff.
### Future considerations
Applying khmer to extremely large data sets with many distinct k-mers requires
a large amount of memory: approximately 446 GB of memory is required to
achieve an false positive rate of 1% for $50\times 10^{9}$ k-mers. It is
possible to reduce the required memory by dividing k-mer space into multiple
partitions and counting k-mers separately for each partition. Partitioning
k-mer space into $M$ partitions results in a linear decrease in the number of
k-mers under consideration, thus reducing the occupancy by a constant factor
$M$ and correspondingly reducing the collision rate. Partitioning k-mer space
is a generalization of the systematic prefix filtering approach, where one
might first count all k-mers starting with AA, then AC, then AG, AT, CA, etc.,
which is equivalent to partitioning k-mer space into 16 equal-sized
partitions. These partitions can be calculated independently, either across
multiple machines or iteratively on a single machine, and the results stored
for later comparison or analysis. This is similar to the approach taken by DSK
[7], and could easily be implemented in khmer.
Further optimization of khmer on single machines, e.g. for multi-core
architectures, is unlikely to achieve significantly greater speed. Past a
certain point k-mer counting is fundamentally I/O bound [23].
Perhaps the most interesting future direction for probabilistic k-mer counting
is that taken by Turtle [9], in which several data structures are provided,
each with different tradeoffs, but with a common API. We hope to pursue this
direction in the future by integrating such approaches into khmer.
## Methods
### Code and data set availability
The version of khmer used to generate the results below is available at
http://github.com/ged-lab/khmer.git, tag ’2013-khmer-counting’. Scripts
specific to this paper are available in the paper repository at
https://github.com/ged-lab/2013-khmer-counting. The IPython[39] notebook file
and data analysis to generate the figures are also available in that github
repository. Complete instructions to reproduce all of the results in this
paper are available in the khmer-counting repository; see README.rst.
### Sequence Data
One human gut metagenome reads data set (MH0001) from the MetaHIT
(Metagenomics of the Human Intestinal Tract) project [40] was used. It
contains approximately 59 million reads, each 44bp long; it was trimmed to
remove low quality sequences.
Five soil metagenomics reads data sets with different size were taken from the
GPGC project for benchmark purpose (see Table 1). These reads are from soil in
Iowa region and they are filtered to make sure there are less than 30% Ns in
the read and each read is longer than 30 bp. The exact data sets used for the
paper are available on Amazon S3 and the instructions to acquire these data
sets are available in the paper repository on github.com.
We also generated four short-read data sets to assess the false positive rate
and miscount distribution. One is a subset of a real metagenomics data set
from the MH0001 data set, above. The second consists of randomly generated
reads. The third and fourth contain reads simulated from a random, 1 Mbp long
genome. The third has a substitution error rate of 3%, and the fourth contains
no errors. The four data sets were chosen to contain identical numbers of
distinct 22-mers. The scripts necessary to regenerate these data are available
in the paper repository on github.com.
### Count-Min Sketch implementation
We implemented the Count-Min Sketch data structure, a simple probabilistic
data structure for counting distinct elements [12]. Our implementation uses
$Z$ independent hash tables, each containing a prime number of counters
$H_{i}$. The hashing function for each hash table is fixed, and reversibly
converts each DNA k-mer (for $k\leq 32$) into a 64-bit number to which the
modulus of the hash table size is applied. This provides $Z$ distinct hash
functions.
To increment the count associated with a k-mer, the counter associated with
the hashed k-mer in each of the $N$ hash tables is incremented. To retrieve
the count associated with a k-mer, the minimum count across all $N$ hash
tables is chosen.
In this scheme, collisions are explicitly not handled, so the count associated
with a k-mer may not be accurate. Because collisions only falsely increment
counts, however, the retrieved count for any given k-mer is guaranteed to be
no less than the correct count. Thus the counting error is one-sided.
### Hash function and khmer implementation
The current khmer hash function works only for $k\leq 32$ and converts DNA
strings exactly into 64-bit numbers. However, any hash function would work.
For example, a cyclic hash would enable khmer to count k-mers larger in size
than 32; this would not change the scaling behavior of khmer at all, and is a
planned extension.
By default khmer counts k-mers in DNA, i.e. strandedness is disregarded by
having the hash function choose the lower numerical value for the exact hash
of both a k-mer and its reverse complement. This behavior is configurable via
a compile-time option.
### Comparing with other k-mer counting programs
We generated k-mer abundance distribution from five soil metagenomic reads
data sets of different sizes using khmer, Tallymer, Jellyfish, DSK, BFCounter,
KMC, Turtle and KAnalyze. If the software does not include function to
generate k-mer abundance distribution directly, we output the frequency of
each k-mer in an output file. We fixed k at 22 unless otherwise noted.
#### khmer:
For khmer, we set hash table sizes to fix the false positive rate at either
1%, 5% or 20%, and used 8 threads in loading the data.
We did the khmer random-access k-mer counting benchmark with a custom-written
Python script khmer-count-kmers which loaded the database file and then used
the Python API to query each k-mer individually.
#### Tallymer:
Tallymer is from the genometools package version 1.3.4. For the suffixerator
subroutine we used: -dna -pl -tis -suf -lcp.
We use the mkindex subroutine to generate k-mer abundance distribution, we
used: -mersize 22.
The Tallymer random access k-mer counting benchmark was done using the
’tallymer search’ routine against both strands; see the script tallymer-
search.sh.
#### Jellyfish:
The Jellyfish version used was 1.1.10 and the multithreading option is set to
8 threads.
Jellyfish uses a hash table to store the k-mers and the size of the hash table
can be modified by the user. When the specified hash table fills up, Jellyfish
writes it to the hard disk and initializes a new hash table. Here we use a
similar strategy as in [6] and chose the minimum size of the hash tables for
Jellyfish so that all k-mers were stored in memory.
We ran Jellyfish with the options as below:
jellyfish count -m 22 -c 2 -C for k=22.
The Jellyfish random access k-mer counting benchmark was performed using the
’query’ routine and querying against both strands; see the script jelly-
search.sh.
#### DSK:
We ran DSK with default parameters with -histo option to generate k-mer
abundance distribution. The DSK version used was 1.5031.
#### BFCounter:
The BFcounter version used was 1.0 and the multithreading option is set to 8
threads.
We ran BFCounter count subroutine with the options as below:
BFCounter count -k 22 -t 8 -c 100000. -n option representing the estimated
number of k-mers is adjusted to the different test data sets.
This subroutine produces the actual count of k-mers in input files.
We ran BFCounter dump subroutine with the options as below: BFCounter dump -k
22.
This subroutine can write k-mer occurrences into a tab-separated text file.
#### KMC:
The KMC version used was 0.3. We ran both kmc and kmc_dump subroutines with
default parameters.
#### Turtle:
The Turtle version used was 0.3. We ran scTurtle32 with the multithreading
option set to 8 threads and -n option representing expected number of frequent
k-mers is adjusted to different test data sets.
#### KAnalyze:
The KAnalyze version used was 0.9.3. We ran count subroutine with default
parameters.
## Acknowledgments
We thank Eric McDonald for technical assistance with optimizing the khmer
codebase.
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## Figure Legends
Figure 1: Comparison of the time it takes for k-mer counting tools to
calculate k-mer abundance histograms, with time (y axis, in seconds) against
data set size (in number of reads, x axis). All programs executed in time
approximately linear with the number of input reads.
Figure 2: Memory usage of k-mer counting tools when calculating k-mer
abundance histograms, with maximum resident program size (y axis, in GB)
plotted against the total number of distinct k-mers in the data set (x axis,
billions of k-mers).
Figure 3: Disk storage usage of different k-mer counting tools to calculate
k-mer abundance histograms in GB (y axis), plotted against the number of
distinct k-mers in the data set (x axis). ∗Note that khmer does not use the
disk during counting or retrieval, although its hash tables can be saved for
reuse.
Figure 4: Time for several k-mer counting tools to retrieve the counts of 9.7m
randomly chosen k-mers (y axis), plotted against the number of distinct k-mers
in the data set being queried (x axis). BFCounter, DSK, Turtle, KAnalyze, and
KMC do not support this functionality.
Figure 5: Relation between average miscount — amount by which the count for
k-mers is incorrect — on the y axis, plotted against false positive rate (x
axis), for five data sets. The five data sets were chosen to have the same
total number of distinct k-mers: one metagenome data set; a set of randomly
generated k-mers; a set of reads, chosen with 3x coverage and 1% error, from a
randomly generated genome; a simulated set of error-free reads (3x) chosen
from a randomly generated genome and a set of E. coli reads.
Figure 6: Relation between percent miscount — amount by which the count for
k-mers is incorrect relative to its true count — on the y axis, plotted
against false positive rate (x axis), for five data sets. The five data sets
are the same as in Figure 5.
Figure 7: Number of unique k-mers (y axis) by starting position within read (x
axis) in an untrimmed E. coli 100-bp Illumina shotgun data set, for k=17 and
k=32. The increasing numbers of unique k-mers are a sign of the increasing
sequencing error towards the 3’ end of reads. Note that there are only 69
starting positions for 32-mers in a 100 base read.
## Tables
Table 1: Benchmark soil metagenome data sets for k-mer counting performance, taken from [11]. Data set | size of file (GB) | number of reads | number of distinct k-mers | total number of k-mers
---|---|---|---|---
subset 1 | 1.90 | 9,744,399 | 561,178,082 | 630,207,985
subset 2 | 2.17 | 19,488,798 | 1,060,354,144 | 1,259,079,821
subset 3 | 3.14 | 29,233,197 | 1,445,923,389 | 1,771,614,378
subset 4 | 4.05 | 38,977,596 | 1,770,589,216 | 2,227,756,662
entire data set | 5.00 | 48,721,995 | 2,121,474,237 | 2,743,130,683
Table 2: Data sets used for analyzing miscounts. Data set | Size of data set file | Number of total k-mers | Number of distinct k-mers
---|---|---|---
Real metagenomics reads | 7.01M | 2,917,200 | 1,944,996
Totally random reads with randomly generated k-mers | 3.53M | 2,250,006 | 1,973,059
Simulated reads from simulated genome with error | 5.92M | 3,757,479 | 2,133,592
Simulated reads from simulated genome without error | 9.07M | 5,714,973 | 1,989,644
Real E. coli reads | 4.85M | 4,004,911 | 2,079,302
Table 3: Iterative low-memory k-mer trimming. The results of trimming reads at unique (erroneous) k-mers from a 5m read E. coli data set (1.4 GB) in under 30 MB of RAM. After each iteration, we measured the total number of distinct k-mers in the data set, the total number of unique (and likely erroneous) k-mers remaining, and the number of unique k-mers present at the 3’ end of reads. | FP rate | bases trimmed | distinct k-mers | unique k-mers | unique k-mers at 3’ end
---|---|---|---|---|---
untrimmed | - | - | 41.6m | 34.1m | 30.4%
khmer iteration 1 | 80.0% | 13.5% | 13.3m | 6.5m | 29.8%
khmer iteration 2 | 40.2% | 1.7% | 7.6m | 909.9k | 12.3%
khmer iteration 3 | 25.4% | 0.3% | 6.8m | 168.1k | 3.1%
khmer iteration 4 | 23.2% | 0.1% | 6.7m | 35.8k | 0.7%
khmer iteration 5 | 22.8% | 0.0% | 6.6m | 7.9k | 0.2%
khmer iteration 6 | 22.7% | 0.0% | 6.6m | 1.9k | 0.0%
filter by FASTX | - | 9.1% | 26.6m | 20.3m | 26.3%
filter by seqtk(default) | - | 8.9% | 17.7m | 12.1m | 12.3%
filter by seqtk(-q 0.01) | - | 15.4% | 9.9m | 5.1m | 5.2%
filter by seqtk(-b 3 -e 5) | - | 8.0% | 34.5m | 27.7m | 25.3%
Table 4: Low-memory digital normalization. The results of digitally normalizing a 5m read E. coli data set (1.4 GB) to C=20 with k=20 under several memory usage/false positive rates. The false positive rate (column 1) is empirically determined. We measured reads remaining, number of “true” k-mers missing from the data at each step, and the number of total k-mers remaining. Note: at high false positive rates, reads are erroneously removed due to inflation of k-mer counts. memory | FP rate | retained reads | retained reads % | true k-mers missing | total k-mers
---|---|---|---|---|---
before diginorm | - | 5,000,000 | 100.0% | 170 | 41.6m
2400 MB | 0.0% | 1,656,518 | 33.0% | 172 | 28.1m
240 MB | 2.8% | 1,655,988 | 33.0% | 172 | 28.1m
120 MB | 18.0% | 1,652,273 | 33.0% | 172 | 28.1m
60 MB | 59.1% | 1,633,182 | 32.0% | 172 | 27.9m
40 MB | 83.2% | 1,602,437 | 32.0% | 172 | 27.6m
20 MB | 98.8% | 1,460,936 | 29.0% | 172 | 25.7m
10 MB | 100.0% | 1,076,958 | 21.0% | 185 | 20.9m
Table 5: E. coli genome assembly after low-memory digital normalization. A comparison of assembling reads digitally normalized with low memory/high false positive rates. The reads were digitally normalized to C=20 (see [21] for more information) and were assembled using Velvet. We measured total length of assembly, as well as percent of true MG1655 genome covered by the assembly using QUAST. memory | FP rate | N contigs | total length(bases) | % of true genome covered
---|---|---|---|---
before diginorm | - | 106 | 4,546,051 | 97.84%
2400 MB | 0.0% | 617 | 4,549,235 | 98.05%
240 MB | 2.8% | 87 | 4,549,253 | 98.04%
120 MB | 18.0% | 86 | 4,549,335 | 98.04%
60 MB | 59.1% | 90 | 4,548,619 | 98.03%
40 MB | 83.2% | 89 | 4,550,599 | 98.11%
20 MB | 98.8% | 85 | 4,550,014 | 98.04%
10 MB | 100.0% | 97 | 4,545,871 | 97.97%
|
arxiv-papers
| 2013-09-11T21:08:59 |
2024-09-04T02:49:50.776380
|
{
"license": "Public Domain",
"authors": "Qingpeng Zhang, Jason Pell, Rosangela Canino-Koning, Adina Chuang\n Howe, C. Titus Brown",
"submitter": "Qingpeng Zhang",
"url": "https://arxiv.org/abs/1309.2975"
}
|
1309.2980
|
# High Fidelity Detection of the Orbital Angular Momentum of Light by Time
Mapping
Paul Bierdz, Minho Kwon, Connor Roncaioli, Hui Deng Department of Physics,
University of Michigan, Ann Arbor, MI 48109, USA [email protected]
###### Abstract
We demonstrate high-fidelity detection of the orbital angular momentum (OAM)
of light using a compact and practical OAM spectrometer that maps the OAM
spectrum to time. The spectrometer consists of a single optical delay loop to
achieve timing mapping, a vortex phase plate that iteratively decreases the
OAM value, and a single mode fibre to distinguish zero from non-zero OAM
states. Light with arbitrarily OAM compositions can be measured. For light
with OAM up to $4\hbar$, we measured an average crosstalk of -21.3 dB, which
is mainly limited by the purity of the input states and optical alignment.
## 1 Introduction
The orbital angular momentum (OAM) of light has been exploited for a wide
range of modern applications. With a discrete, yet unbounded Hilbert space
[1], OAM has been used to increase the capacity of free-space communications
[2, 3], generate high order entanglement [4, 5] and increase security of non-
classical communications [6]. The unique topological properties of OAM states
have been utilized to create optical tweezers that can apply varying degrees
of torque [7, 8], high-sensitivity vortex coronagraphs for extrasolar planet
detection[9] and revealing topological properties of objects[10]. A recent
paper suggests that light emitted nearby rotating black holes is twisted such
that the angular momentum of the black hole can be inferred from the OAM
spectrum [11].
For these applications, it is often crucial to discern different OAM states
with high fidelity. Several methods have been demonstrated for measuring the
OAM spectrum including those using fork diffraction gratings coupled with
single mode fibres (SMF) [4], cascading Mach-Zehnder interferometers [12] and
transformation optics [13]. However, all these methods map different OAM
states to spatially separated modes; as a result, complexity and size of the
experimental setups increase with the highest measurable OAM state. Fidelity
of the detection is often limited by the complexity of the experimental setup
and optical diffraction. There are other techniques that do not spatially
separate the modes. Courtial et. al. proposed a scheme based on OAM-dependent
rotational Doppler shift [14]; yet experimental implementation has not been
possible for optical wavelengths. Recently, we proposed a scheme based on OAM
to time mapping [15]. The scheme utilizes the concept of quantum
counterfactual measurement and quantum Zeno effect to allow non-destructive
sorting of the OAM spectrum with unit efficiency, if optical loss or
misalignment is neglected. Experimental implementation with lossy optics,
however, will significantly reduce the efficiency for high-order OAM state.
Here we demonstrate a simplified, and practical, OAM-to-time mapping scheme,
which achieved a record high extinction ratio among five OAM states at an
operation speed of 80 MHz. Since OAM is mapped to time, the same apparatus can
be used to measure an arbitrary number of unique OAM states. We note that a
similar scheme was also adopted in a recent experiment on time-division
multiplexing of OAM, although the fidelity and repetition rate were orders of
magnitude lower than reported here [16]. We also note that, without employing
a quantum Zeno investigator, time-division schemes are not suitable for
improving channel capacity in communication.
## 2 Principle of the OAM spectrometer
As shown in figure 1a, our OAM spectrometer consists of a single optical loop
to perform time mapping, a vortex phase plate (VPP) of topological charge 1 as
an OAM ladder operator, and an SMF as a filter for the fundamental Gaussian
mode with zero OAM. Consider an incident pulse consisting of a fraction
$\beta_{\ell}$ of OAM components with OAM value $\ell$, where $\ell$ is an
integer and $\sum_{\ell=-\infty}^{\infty}\beta_{\ell}=1$. The optical loop
converts the pulse into a sequence of pulses equally separated by the round
trip propagation time $T$ in the loop. The loop needs an even number of
reflections to maintain the same sign of OAM. Per loop, the VPP decreases the
OAM value of each OAM component by $1$. Hence, after $N$ loops, the fractional
$\beta_{\ell}$ of the original pulse will have OAM value $\ell-N$. Only the
component with zero OAM value in each pulse can pass through the SMF to be
detected. Thus, the OAM component $\beta_{\ell}$ with OAM value $\ell$ in the
original pulse will exit the spectrometer at a time $t=lT$ (figure 1b).
Figure 1: (a) Schematic of the OAM spectrometer, consisting of an optical loop
that converts an input pulse into a sequence of pulses equally spaced in time,
a VPP that decreases the OAM value by 1 per pass, and an SMF to filter out
states with non-zero OAM. (b) Showing how the OAM value changes in time for an
example input state with $\ell_{0}=3$.
### 2.1 Energy distribution
The distribution of the input-pulse energy among the sequence of pulses can be
pre-calibrated and controlled by the beam splitter with polarization optics.
We control the distribution with a polarizing beam-splitter (PBS) and half-
wave plate (labelled as R in figure 1a) in the loop. The incident light is
first set to be linearly polarized at an angle $\phi_{0}$ with respect to the
vertical. The PBS splits off the fraction of $\sin^{2}\\!\phi_{0}$ into the
first time window and sends the rest, now vertically polarized, into the loop.
The wave plate R, rotated at an angle $\theta/2$ with respect to the vertical,
rotates the vertically polarized light by an angle $\theta$. The corresponding
energy distribution is described below:
$\displaystyle P(l)=\cases{\sin^{2}\\!\phi_{0}&\mbox{if $\ell=0$,}\\\
\cos^{2}\\!\phi_{0}\cos^{2}\\!\theta&\mbox{if $\ell=1$,}\\\
\cos^{2}\\!\phi_{0}\sin^{4}\\!\theta\cos^{2(\ell-2)}\\!\theta&\mbox{if
$\ell\geq 2$.}}$ (1)
Here $\ell$ represents the $\ell$th output pulse. Alternatively, the PBS and
wave plate can be replaced by a non-polarizing beam-splitter with a chosen
splitting ratio for polarization-insensitive measurements. This would allow
for information to be encoded in the polarization degree of freedom.
## 3 Experimental implementation
We implemented an OAM spectrometer as illustrated in figure 1(a). To test its
performance, we used OAM eigenstates as input pulses and detected the output
pulses using a Hamamatsu streak camera.
The input pulse was generated by diffracting a pulsed Gaussian laser beam off
fork-diffraction patterns [17] on a Holoeye PLUTO LCoS spatial light modulator
(SLM) with a period of about 10 lines per millimeter and pixel pitch of
8$\mu$m. The initial laser beam was from a Tsunami Ti-Sapphire laser centered
at 730 nm, with a pulse width of 100 fs and repetition rate of 80 MHz, and
collimated with a 10x objective lens from an SMF ensuring an initial $M^{2}$
close to 1 and a spot size within the frame of the SLM.
We first verify the input states by measuring its far-field intensity and
phase distributions, as shown in figure 2. The intensity distributions were
measured with a charge-coupled device. The higher-order OAM beams show a
larger spatial size as expected. Inhomogeneity in the radial intensity
distribution reflects the lack of purity of the OAM value.
We measured the phase distribution by interfering the OAM beam with a
reference Gaussian beam from the original laser. This produces a forked
interference pattern. From this interference pattern we extract the phase
front, as plotted in figure 2(f-j), by performing a Fourier transform,
selecting the first diffracted mode, centring and then applying an inverse
Fourier transform. The phase should twist about by $2\pi\ell$ for an
$\ell^{\mbox{th}}$-order OAM beam. Deviation from this results implies impure
initial OAM states from optical aberrations or phase ripples from modulating
the liquid crystals.
To calibrate the time-mapping and energy distribution of the OAM spectrometer,
we used the zero-OAM Gaussian state as the input and replaced the VPP by a
flat glass plate of the same thickness. The output from the SMF was detected
by a Hamamatsu streak camera with a time resolution of 0.02 ns. As shown in
figure 1(b), the input pulse was mapped into a sequence of pulses separated by
$T=1.03$ ns. The timing and dynamic range of the streak camera allowed us to
measure the first five output pulses, or, the first five OAM states, at the
laser repetition rate of 80 MHz. The experiment can also be redesigned to
allow the measurement of a larger number of pulses by using different laser
repetition rates, loop sizes or measurement devices.
Figure 2: (a-e) The intensity measured with a camera of the different initial
OAM beams from the SLM with $\ell=0-4$. (f-j) The phase fronts of the input
OAM beams. They are calculated from the interference patterns between the OAM
beams and a reference Gaussian beam, as explained in the text.
The energy distribution among the pulses was calculated by integrating the
intensity over each output pulse. The results are shown in figure 3(a) (the
$\ell_{0}=0$ curve). We performed the same measurement with the input state
$\ell_{0}=1$. The ratio of the two energy normalization curves at each $\ell
T$ gives an estimate of the misalignment of the SMF-coupling (figure 3(b)) for
the $\ell^{\mbox{th}}$ output pulse, as will be discussed in detail in Section
4.2 later. The uniformity of ratio among different $\ell$ confirms that the
loop was well aligned.
Finally, we benchmark the performance of the OAM spectrometer by comparing the
correct versus incorrect detection of an input OAM eigenstate. Figure 4(a-e)
show examples of the output pulse sequences for input states $\ell_{0}=0-4$.
The integrated intensity under each peak versus the input state $\ell_{0}$ and
the output time bin $\ell$ is shown in figure 4(f), where each row has been
renormalized to the peak of the correct detection. We define the crosstalk to
be the ratios between the left and right adjacent incorrect detections versus
the correct detection, i.e., $\frac{P(\ell_{0}=\ell,N=\ell\pm
1)}{P(\ell_{0}=\ell,N=\ell)}$ where $P(\ell_{0},N)$ refers to the intensity
measured at time $NT$ with input OAM state $\ell_{0}$. For $\ell_{0}=0\mbox{
to }4$ we measured crosstalk values: $-12.3$, $-18.8$, $-20.0$, $-21.1$,
$-21.6$, $-23.2$, $-24.8$ and $-28.2$ dB. The geometric mean of the crosstalk
is $7.47\times 10^{-3}$ or $-21.3$ dB.
Figure 3: (a) Unnormalized power through the fibre versus the number of passes
through the loop without a VPP and different initial OAM values
($\ell_{0}=0,1$) set by the SLM. (b) Computed misalignment of SMF.
Figure 4: (a-e) Streak camera images for OAM eigenstates ($\ell=0,2,4$)
normalized based on (1). (f) Tabulated results of the OAM spectrometer up to
$\ell=4$.
## 4 Analysis of the spectrometer performance
Below we simulate the performance of the spectrometer and analyze the main
sources of error. We model the laser pulses as a superposition of Laguerre-
Gaussian (LG) modes, which are a complete orthonormal set of solutions to the
paraxial wave equation. Each optical element operating on the pulse maps each
LG mode into a different superposition of other LG modes. In particular, we
focus on SMF, VPP and SLM. The remaining optics: mirrors, beam splitters,
waveplates, etc. are insensitive to the transverse mode of the beam and thus
preserve the LG mode. We then simulate the propagation of an input pulse
through the spectrometer, and the coefficients of superposition obtained after
the SMF corresponds to the measurement results.
### 4.1 Laguerre-Gaussian modes
The electric fields of LG modes can be described in cylindrical coordinates
$(\rho,\phi,z)$ as [18]:
$\displaystyle
u_{\ell,p}(\rho,\phi,z)=\langle\rho,\phi,z|\ell,p\rangle=\frac{1}{w(z)}\sqrt{\frac{2}{\pi}\frac{p!}{(p+|\ell|)!}}\left(\frac{\sqrt{2}\rho}{w(z)}\right)^{|\ell|}\exp\left({-\frac{\rho^{2}}{w^{2}(z)}}\right)L_{p}^{|\ell|}\left(\frac{2\rho^{2}}{w^{2}(z)}\right)\rme^{\rmi\ell\phi}$
$\displaystyle\exp\left(\frac{\rmi k_{0}\rho^{2}z}{2\left(z^{2}+z_{\rm
R}^{2}\right)}\right)\exp\left({-\rmi(2p+|\ell|+1)\tan^{-1}\left(\frac{z}{z_{\rm
R}}\right)}\right)$ (2)
Here $\ell\hbar$ ($\ell\in\mathbb{Z}$) is the OAM per photon, $p\geq 0$ labels
the radial modes, $w(z)$ is the beam waist, $z_{\rm R}$ is Rayleigh range, and
$k_{0}$ is the wave number of the fundamental Gaussian mode.
### 4.2 Single mode fibre
Ideally, an SMF selects only the fundamental Gaussian mode ($\ell=0,p=0$)
while all the other orthogonal spatial modes do not propagate through. In
reality, higher order LG modes may still couple through the fibre due to
finite apertures of optics and the fibre, imperfections of the fibre,
mismatched beam waists between free-space and fibre-modes, and transverse
misalignment between the propagation axes of the free-space and fibre-modes.
In our experiments, the last effect, the transverse misalignment, is by far
the dominant, while the other effects were too small to be measured. Hence we
restrict our discussion to the transverse misalignment only. For a
misalignment of $\Delta$ (in units of beam waist) between the two optical
axes, the coupling efficiency of the LG beam $(\ell,p)$ through the SMF is
given by (A.2):
$|{}_{\Delta}\langle
0,0|\ell,p\rangle|^{2}=\frac{1}{p!(p+|\ell|)!}\left(\frac{\Delta^{2}}{2}\right)^{2p+|\ell|}\rme^{-\Delta^{2}}$
(3)
Experimentally, the above misalignment of the SMF can be estimated using the
data for calibrating the energy distribution, as shown in figure 3(a). Nonzero
output intensity from an input pulse of $\ell_{0}=1$ modes implies a
misaligned (or imperfect) SMF. The amount of misalignment $\Delta$ can be
calculated from the ratio of the two output intensities corresponding to the
two input modes $\ell_{0}=0$ versus $\ell_{0}=1$ via:
$\frac{P(\ell_{0}=1)}{P(\ell_{0}=0)}=\frac{|{}_{\Delta}\langle
0,0|M_{\mbox{VPP1}}|0,0\rangle|^{2}}{|{}_{\Delta}\langle 0,0|0,0\rangle|^{2}}$
(4) $|{}_{\Delta}\langle
0,0|M_{\mbox{VPP1}}|0,0\rangle|^{2}=\frac{\pi}{8}\Delta^{2}\rme^{-\frac{3}{2}\Delta^{2}}\left(I_{0}\left(\frac{\Delta^{2}}{4}\right)+I_{1}\left(\frac{\Delta^{2}}{4}\right)\right)^{2}$
(5)
The left hand side of (4) is measured experimentally. $M_{\mbox{VPP1}}$
represents the operation by the SLM to create the input beam with $\ell_{0}=1$
and in a superposition of different $p$-states [19]. Equation (5) is obtained
using (3) and (6) (See A.3 for details). We plotted the derived misalignments
in figure 3(b). The misalignment is less than 10.1% for all loops and
converges to 3.7% for higher loop numbers, indicating only a minor
accumulative loop misalignment.
### 4.3 VPP/SLM
The VPP and SLM are used to change the OAM state. The SLM, when applying the
forked phase pattern, is mathematically equivalent to VPP for small
diffraction angles [20, 21]. For a VPP of topological charge $\beta$, its
operation on the laser beam can be described by a four-dimensional tensor
$M_{VPP\beta}(z)$. We solve for the tensor elements
$m_{\ell_{1},p_{1};\ell_{2},p_{2};\beta}(z)$ analytically (A.1) in the LG-mode
basis:
$\displaystyle
m_{\ell_{1},p_{1};\ell_{2},p_{2};\beta}(z)=\left[\sqrt{\frac{p_{1}!p_{2}!}{(p_{1}+|\ell_{1}|)!(p_{2}+|\ell_{2}|)!}}\right.\sum_{k=0}^{p_{1}}\sum_{m=0}^{p_{2}}(-1)^{k+m}\frac{(p_{1}+|\ell_{1}|)!}{(p_{1}-k)!(|\ell_{1}|+k)!k!}$
$\displaystyle\frac{(p_{2}+|\ell_{2}|)!}{(p_{2}-m)!(|\ell_{2}|+m)!m!}\left.\Gamma\left(\frac{|\ell_{1}|+|\ell_{2}|}{2}+k+m+1\right)\right]$
$\displaystyle\left[\exp\left(\rmi(2(p_{1}-p_{2})+|\ell_{1}|-|\ell_{2}|)\tan^{-1}\left(\frac{z}{z_{\rm
R}}\right)\right)\right]$
$\displaystyle\left[\frac{\exp[2\pi\rmi(\ell_{2}+\beta-\ell_{1})]-1}{2\pi\rmi(\ell_{2}+\beta-\ell_{1})}\right]$
(6)
This equation consists of a product of three terms denoted by square brackets.
The first term with the double sum comes from the amplitude overlap between
the different LG modes. The second term, consisting of an exponential,
describes the effects of Gouy phases. The third term comes from OAM
conservation and can describe the error in the topological charge of the VPP.
We discuss the effect of each term below.
#### 4.3.1 Gouy phase effects
Different LG modes have different Gouy phases which also vary differently as
the beam propagates; therefore, whenever there is a superposition of modes,
there can be interference effects between the modes that varies under free
space propagation [22] and can lead to the loss of efficiency. However, the
effect does not produce crosstalk between different $\ell$-states, as OAM
remain unchanged. In our experiment, the effect of the Gouy phase accounts for
less than $0.05\%$ loss in efficiency. This upper bound of $0.05\%$ was
calculated based on distances between the optics in the experimental setup and
using (6). In general, the effect is negligible when the propagation distance
is much smaller than the Rayleigh range ($z\ll z_{\rm R}$). When the size of
the OAM spectrometer becomes comparable to the Rayleigh range, the use of 4-f
systems between all phase elements would eliminate any Gouy phase effects.
#### 4.3.2 Topological charge error
If the topological charge of the VPP is an integer, then the OAM of the beam
is changed by that amount. If the laser’s wavelength is different from the
nominal wavelength of the VPP, or if the beam comes at a skew angle to the
VPP, then the VPP will appear thicker or thinner. As a result, the OAM of the
beam will be changed by a fractional value instead, or equivalently, become a
superposition of LG modes of many OAM values. This leads to a loss in both the
efficiency and fidelity of the OAM spectrometer.
Such topological charge error can be modelled by non-integer $\beta$ in (6).
In our experiment, both the laser wavelength and the angle of incidence are
tightly controlled. Even if we consider a rather generous 0.5% error in the
topological charge, it results in less than 0.05% loss of efficiency and less
than -34dB crosstalk. Therefore the error in the topological charge is
negligible.
#### 4.3.3 Lateral misalignment
Although topological charge error will reduce the fidelity, it is, by far, not
the leading cause. If the optical axes of the beam and VPP are displaced
relative to each other, the VPP produces a superposition of not just $p$, but
$\ell$ states [23] and thus reduces the fidelity. Such lateral misalignment of
the VPP affects all three terms in (6). To model it, we first express the VPP
tensor $M_{\beta}=\rme^{\rmi\beta\phi}$ in a coordinate displaced from the
common optical axis of the spectrometer. We then numerically evaluate
$m_{\ell_{1},p_{1};\ell_{2},p_{2};\beta}(z)$ in a subspace of
${\ell,p}=[-7,7]\times[0,10]$. This chosen subspace yields less than 2% error
on theoretical crosstalk.
To compare with the experimental result, we calculated the average crosstalk
versus VPP misalignment, assuming a transverse misalignment of the SMF of 3.7%
to 10.1% as obtained in 4.2. The results are shown in 5(a), where the
x-direction corresponds to the direction of SMF misalignment. The measured
crosstalk of $-21.3$ dB (3) corresponds to a VPP misalignment of 4.0% to 6.1%
of the beam waist $w_{0}$. A larger (smaller) VPP misalignment is allowed when
its direction is aligned with (perpendicular to) that of the SMF misalignment.
### 4.4 Limiting factors of fidelity
In short, we show that misalignments of 3.7% to 10.1% at the SMF and 4.0% to
6.1% at the VPP are the leading causes for the loss of fidelity in our
experiment, with an average crosstalk of -21.3 dB. Gouy phase effects do not
affect fidelity and topological charge error contributes less than -34 dB to
crosstalk.
The fidelity of the spectrometer can be further increased by improving the
optical alignment at the VPP and SMF. We show in figure 5(b) the average
crosstalk as a function of the VPP misalignment alone, neglecting SMF
misalignment, and as a function of the SMF misalignment alone, neglecting VPP
misalignment. It is more sensitive to VPP misalignment than to SMF
misalignment, but decreases superlinearly with the reduction of either
misalignment. With a VPP (SMF) misalignment of $<1\%$, the crosstalk is
reduced to $<-34.4$ dB ($<-43.4$ dB). When including both SMF and VPP
misalignments, the crosstalk varies depending on the relative angle between
the directions of the two misalignments. The crosstalk is the greatest when
the misalignments are in perpendicular directions from each other. We show
this upper bound crosstalk in 5(c). When the misalignments are in opposite
directions, we generally see a reduction in crosstalk, sometimes up to 20 dB.
Curiously, if the SMF is misaligned, crosstalk can be _decreased_ by
introducing some VPP misalignment, and vice versa. This is because with
multiple elements misaligned, they can compensate each other. However, zero
crosstalk is reached only when both misalignments are exactly zero.
Figure 5: Calculated the geometric mean of crosstalk (in dB) as a function of
the lateral misalignment of the SMF and VPP (normalized by $w_{0}$). (a)
Crosstalk (in dB) as a function of the lateral misalignment of the VPP. SMF
misalignment (measured in 3(c)) is taken along the x direction for
convenience, since only the relative angle between VPP and SMF misalignment is
important. The measured crosstalk of $-20.9$ to $-21.6$ dB corresponds to the
values between the two black lines. (b) The crosstalk versus VPP misalignment
with no SMF misalignment (green solid line), and the crosstalk versus SMF
misalignment with no VPP misalignment (blue dashed line). (c) The upper bound
of crosstalk (in dB) as a function of the SMF and VPP misalignments. The upper
bound corresponds to when the directions of the SMF and VPP misalignments are
perpendicular to each other. The crosstalk can be significantly reduced if the
directions of the SMF and VPP lateral displacements are anti-parallel to each
other.
As misalignment is reduced, other subtle effects will need to be taken under
consideration. Such as imperfections in the surface roughness or other
aberrations in the optics, which could further reduce efficiency and fidelity.
Eventually, the fidelity will be limited by overlap in time between adjacent
pulses. We did not include these effects in our model.
## 5 Conclusions
In summary, we present a practical OAM spectrometer that can map OAM to time
up to arbitrarily large OAM values. We have shown an average nearest neighbor
crosstalk of -21.3dB among 5 OAM states, limited mainly by optical
misalignment. The high fidelity of the demonstrated OAM spectrometer may
enable accurate measurements of topological properties of objects such as in
spiral imaging [10] and the angular momentum of black holes, which is encoded
into the OAM spectrum of light from the accretion disc due to strong
gravitational effects predicted by general relativity [11]. We also
demonstrated speeds (80 MHz) orders of magnitude faster than the switching
times of SLMs. Miniaturization of optics could allow for GHz detection rate.
## Acknowledgments
We would like to thank Professor Miles Padgett for his helpful advice. And
also thank Professor Theo Lasser and Matthias Geissbühler for their
_Morgenstemning_ colormap [24].
## Appendix A
In this appendix section we seek to derive the various actions and overlaps
between the LG modes and optics. We first derive the action of the VPP in the
LG basis. Then we derive the overlap of any LG mode with a misaligned SMF.
Lastly, we derive the overlap of OAM states created by a VPP with a misaligned
SMF. This is done by explicit integration and definitions of various special
functions. We use the inner product between Laguerre-Gaussian modes defined in
2 in cylindrical coordinates:
$\displaystyle\langle\ell_{1},p_{1}|\ell_{2},p_{2}\rangle=\int_{0}^{\infty}\rho\rmd\rho\int_{0}^{2\pi}\rmd\phi\langle\ell_{1},p_{1}|\rho,\phi,z\rangle\langle\rho,\phi,z|\ell_{2},p_{2}\rangle$
$\displaystyle\times\int_{0}^{\infty}\rho\rmd\rho\int_{0}^{2\pi}\rmd\phi
u_{\ell_{1},p_{1}}^{*}(\rho,\phi,z)u_{\ell_{2},p_{2}}(\rho,\phi,z)$ (7)
### A.1 Derivation of VPP tensor
We derive the four dimensional tensor of the action of a VPP in the LG basis
as seen in (6). We start by explicitly writing down the overlap integral:
$\langle\ell_{1},p_{1}|M_{\mbox{VPP}\beta}|\ell_{2},p_{2}\rangle=\int_{0}^{\infty}\rho\rmd\rho\int_{0}^{2\pi}\rmd\phi
LG^{*}_{\ell_{1},p_{1}}(\rho,\phi)\rme^{+\rmi\beta\phi}LG_{\ell_{2},p_{2}}(\rho,\phi)$
(8)
The LG functions are given by (2). Due to separability of variables, we will
solve the $\phi$-integral first:
$\int_{0}^{2\pi}\rmd\phi\rme^{-\rmi\ell_{1}\phi}\rme^{+\rmi\beta\phi}\rme^{+\rmi\ell_{2}\phi}=\frac{\exp\left(2\pi\rmi\left(\ell_{2}+\beta-\ell_{1}\right)\right)-1}{\rmi\left(\ell_{2}+\beta-\ell_{1}\right)}$
(9)
In the limit where $\beta$ is an integer, the $\phi$-integral yields the OAM
conserving solution of $2\pi\delta_{\ell_{1}+\beta-\ell_{2},0}$, where
$\delta_{a,b}$ is the Kronecker delta. The remaining $\rho$-integral is found
by expanding the rest of the LG modes and the generalized Laguerre polynomials
given by (10). This yields a finite polynomial of degree
$|\ell_{1}|+|\ell_{2}|+2p_{1}+2p_{2}+1$ (the +1 is from the Jacobian)
multiplied by the Gaussian function
$\exp\left(\frac{-\rho^{2}}{2w_{0}^{2}}\right)$. This integral can be easily
solved by using one of the definitions of the Gamma function as shown in (11)
to yield the result in the text (6).
$L_{p}^{k}(x)=\sum_{m=0}^{p}(-1)^{m}\frac{(p+k)!}{(p-m)!(k+m)!m!}x^{m}$ (10)
$\int_{0}^{\infty}\rmd
xx^{k}\rme^{-x^{2}}=\frac{1}{2}\Gamma\left(\frac{1+k}{2}\right)$ (11)
### A.2 Overlap between fundamental Gaussian and Laguerre-Gaussian
We first derive the overlap between a misaligned fundamental Gaussian mode and
generic higher order LG modes. This is (3) in the text. We start with the
fundamental Gaussian mode defined in cylindrical coordinates ($\rho,\phi,z$)
with an origin shifted by $w_{0}\Delta$ along the $\phi=0$ direction. The
propagation direction is along the z-axis.
${}_{\Delta}\langle
0,0|\rho,\phi,z\rangle=\frac{1}{w_{0}}\sqrt{\frac{2}{\pi}}\exp\left({\frac{-\left(\rho^{2}+2\rho
w_{0}\Delta\cos\phi+w_{0}^{2}\Delta^{2}\right)}{w_{0}^{2}}}\right)$ (12)
If the origin were shifted along a different $\phi$ angle, the only change
would be the overall phase which does not affect the intensity. These
calculations are performed at the beam waist since the beam should be focused
on the fibre to have any substantial coupling. With a change of variables
$x=\frac{\sqrt{2}\rho}{w_{0}}$ the complete overlap integral is now:
$\displaystyle{}_{\Delta}\langle
0,0|\ell,p\rangle=\int_{0}^{\infty}\int_{0}^{2\pi}x\rmd
x\rmd\phi\frac{1}{\pi}\sqrt{\frac{p!}{(p+|\ell|)!}}x^{|\ell|}\exp\left(\frac{-x^{2}}{2}\right)L_{p}^{|\ell|}\left(x^{2}\right)\rme^{\rmi\ell\phi}$
$\displaystyle\exp\left(-\left(\frac{x^{2}}{2}+\sqrt{2}x\Delta\cos\phi+\Delta^{2}\right)\right)$
(13)
The $\phi$-integral is related to the well known Bessel integral (14), so
(A.2) becomes (15).
$J_{n}(z)=\frac{1}{2\pi\rmi^{n}}\int_{0}^{2\pi}\rmd\phi\rme^{\rmi
z\cos\phi}\rme^{\rmi n\phi}$ (14)
${}_{\Delta}\langle
0,0|\ell,p\rangle=\sqrt{\frac{p!}{(p+|\ell|)!}}\int_{0}^{\infty}x\rmd
xx^{|\ell|}\rme^{-\left(x^{2}+\Delta^{2}\right)}L_{p}^{|\ell|}\left(x^{2}\right)2\rmi^{-|\ell|}(-1)^{|\ell|}J_{|\ell|}\left(\rmi\sqrt{2}\Delta
x\right)$ (15)
The Bessel function can be converted into an infinite sum of generalized
Laguerre polynomials as given in (16) below with
$\sqrt{w}=\rmi\Delta/\sqrt{2}$. Using (16) and the orthogonality relationship
between generalized Laguerre polynomials (17). Equation (15) magnitude squared
surprisingly reduces to the rather simple expression of (3) in the text.
$\displaystyle\sum_{n=0}^{\infty}\frac{L_{n}^{k}(x)}{\Gamma(n+k+1)}w^{n}=e^{w}(xw)^{\frac{-k}{2}}J_{k}\left(2\sqrt{xw}\right)$
(16)
$\displaystyle\int_{0}^{\infty}\rmd
x\rme^{-x}L_{n}^{k}(x)L_{m}^{k}(x)=\frac{(n+k)!}{n!}\delta_{m,n}$ (17)
### A.3 Overlap between fundamental Gaussian and vortex phase plate mode
We create states with OAM by passing a Gaussian beam through a VPP. This state
can be expressed as a sum of Laguerre Gaussian modes as given in (18), which
is a special case of the more general equation solved in section A.1. In the
previous section, we calculated the overlap of a shifted fundamental Gaussian
mode and any Laguerre Gaussian mode. Therefore, by combining these two
calculations, we can derive the overlap of a shifted Gaussian mode and an OAM
mode created by a VPP (19), here $z=\frac{\Delta^{2}}{2}$.
$\displaystyle
M_{\mbox{VPPN}}|0,0\rangle=\sum_{p=0}^{\infty}m_{0,0;N,p;N}|N,p\rangle$ (18)
$\displaystyle\mbox{where
}m_{0,0;N,p;N}=\frac{N}{2}\sqrt{\frac{1}{(p+N)!}}\Gamma\left(p+\frac{N}{2}\right)$
$\displaystyle{}_{\Delta}\langle
0,0|M_{\mbox{VPPN}}|0,0\rangle=\frac{e^{-z}z^{\frac{N}{2}}(-1)^{N}N}{2}\sum_{p=0}^{\infty}\frac{(-z)^{p}\Gamma\left(p+\frac{N}{2}\right)}{(p+N)!p!}$
(19)
Equation (19) can be written in the form of a regularized confluent
hypergeometric function (20) which can be evaluated using
(21)111http://functions.wolfram.com/07.21.03.0011.01 to yield the result in
(22), where $I_{\alpha}(x)$ is the modified Bessel function of the first kind.
The expression can be further simplified by noting
$I_{\alpha}(-z)=\rmi^{2\alpha}I_{\alpha}(z)$. In the case of $N=1$ and taking
the magnitude squared, this reduces to the form in the text (5).
${}_{1}\tilde{F}_{1}(a;b;z)=\sum_{k=0}^{\infty}\frac{\Gamma(a+k)}{\Gamma(b+k)\Gamma(a)}\frac{z^{k}}{k!}$
(20)
${}_{1}\tilde{F}_{1}\left(\frac{N}{2};N+1;z\right)=\frac{e^{\frac{z}{2}}\sqrt{\pi}}{2\Gamma\left(\frac{N}{2}+1\right)}z^{\frac{1}{2}-\frac{N}{2}}\left(I_{\frac{N}{2}-\frac{1}{2}}\left(\frac{z}{2}\right)-I_{\frac{N}{2}+\frac{1}{2}}\left(\frac{z}{2}\right)\right)$
(21)
$\frac{\rme^{\frac{-3\Delta^{2}}{4}}\Delta^{2}\sqrt{\pi}}{2\sqrt{2}}\left(I_{\frac{N}{2}-\frac{1}{2}}\left(\frac{\Delta^{2}}{4}\right)+I_{\frac{N}{2}+\frac{1}{2}}\left(\frac{\Delta^{2}}{4}\right)\right)$
(22)
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* [8] David G. Grier. A revolution in optical manipulation. Nature, 424(6950):810–816, August 2003. URL: http://www.nature.com/nature/journal/v424/n6950/full/nature01935.html, doi:10.1038/nature01935.
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* [11] Fabrizio Tamburini, Bo Thidé, Gabriel Molina-Terriza, and Gabriele Anzolin. Twisting of light around rotating black holes. Nature Physics, 7(3):195–197, 2011. URL: http://www.nature.com/nphys/journal/v7/n3/full/nphys1907.html, doi:10.1038/nphys1907.
* [12] Jonathan Leach, Miles J. Padgett, Stephen M. Barnett, Sonja Franke-Arnold, and Johannes Courtial. Measuring the orbital angular momentum of a single photon. Physical Review Letters, 88(25):257901, June 2002. URL: http://link.aps.org/doi/10.1103/PhysRevLett.88.257901, doi:10.1103/PhysRevLett.88.257901.
* [13] Gregorius C. G. Berkhout, Martin P. J. Lavery, Johannes Courtial, Marco W. Beijersbergen, and Miles J. Padgett. Efficient sorting of orbital angular momentum states of light. Physical Review Letters, 105(15):153601, October 2010. URL: http://link.aps.org/doi/10.1103/PhysRevLett.105.153601, doi:10.1103/PhysRevLett.105.153601.
* [14] J. Courtial, K. Dholakia, D. A. Robertson, L. Allen, and M. J. Padgett. Measurement of the rotational frequency shift imparted to a rotating light beam possessing orbital angular momentum. Physical Review Letters, 80(15):3217–3219, April 1998. URL: http://link.aps.org/doi/10.1103/PhysRevLett.80.3217, doi:10.1103/PhysRevLett.80.3217.
* [15] Paul Bierdz and Hui Deng. A compact orbital angular momentum spectrometer using quantum zeno interrogation. Optics Express, 19(12):11615–11622, June 2011. URL: http://www.opticsexpress.org/abstract.cfm?URI=oe-19-12-11615, doi:10.1364/OE.19.011615.
* [16] Ebrahim Karimi, Lorenzo Marrucci, Corrado de Lisio, and Enrico Santamato. Time-division multiplexing of the orbital angular momentum of light. Optics Letters, 37(2):127–129, January 2012. URL: http://ol.osa.org/abstract.cfm?URI=ol-37-2-127, doi:10.1364/OL.37.000127.
* [17] VY Bazhekov, MV Vasnetsov, and M S Soskin. Laser-beams with screw dislocation in their wave-fronts. JETP Letters, 52(8):429–431, October 1990.
* [18] Alison M. Yao and Miles J. Padgett. Orbital angular momentum: origins, behavior and applications. Advances in Optics and Photonics, 3(2):161–204, June 2011. URL: http://aop.osa.org/abstract.cfm?URI=aop-3-2-161, doi:10.1364/AOP.3.000161.
* [19] Ebrahim Karimi, Gianluigi Zito, Bruno Piccirillo, Lorenzo Marrucci, and Enrico Santamato. Hypergeometric-Gaussian modes. Optics Letters, 32(21):3053–3055, November 2007. URL: http://ol.osa.org/abstract.cfm?URI=ol-32-21-3053, doi:10.1364/OL.32.003053.
* [20] N. R. Heckenberg, R. McDuff, C. P. Smith, H. Rubinsztein-Dunlop, and M. J. Wegener. Laser beams with phase singularities. Optical and Quantum Electronics, 24(9):S951–S962, September 1992\. URL: http://link.springer.com/article/10.1007/BF01588597, doi:10.1007/BF01588597.
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* [22] J. Arlt and M. J. Padgett. Generation of a beam with a dark focus surrounded by regions of higher intensity: the optical bottle beam. Optics Letters, 25(4):191–193, February 2000. URL: http://ol.osa.org/abstract.cfm?URI=ol-25-4-191, doi:10.1364/OL.25.000191.
* [23] Gabriel Molina-Terriza, Juan P. Torres, and Lluis Torner. Management of the angular momentum of light: Preparation of photons in multidimensional vector states of angular momentum. Physical Review Letters, 88(1):013601, December 2001. URL: http://link.aps.org/doi/10.1103/PhysRevLett.88.013601, doi:10.1103/PhysRevLett.88.013601.
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|
arxiv-papers
| 2013-09-11T21:56:31 |
2024-09-04T02:49:50.787967
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Paul Bierdz, Minho Kwon, Connor Roncaioli, Hui Deng",
"submitter": "Paul Bierdz",
"url": "https://arxiv.org/abs/1309.2980"
}
|
1309.2988
|
EPJ Web of Conferences INPC 2013
11institutetext: Mississippi State University, MS 39762-5167, USA
# Feasibility of the Spin-Light Polarimetry Technique for Longitudinally
Polarized Electron Beams
Prajwal Mohanmurthy 11 [email protected] Dipangkar Dutta 11
###### Abstract
A novel polarimeter based on the asymmetry in the spacial distribution of
synchrotron radiation will make for a fine addition to the existing Møller and
Compton polarimeters. The spin light polarimeter consists of a set of wiggler
magnet along the beam that generate synchrotron radiation. The spacial
distribution of synchrotron radiation will be measured by ionization chambers.
The up-down (below and above the wiggle) spacial asymmetry in the transverse
plain is used to quantify the polarization of the beam. As a part of the
design process, effects of a realistic wiggler magnetic field and an extended
beam size were studied. The perturbation introduced by these effects was found
to be negligible. Lastly, a full fledged GEANT-4 simulation was built to study
the response of the ionization chamber.
## 1 Introduction
A 1993 proposal from Karabekov and Rossmanith explored the possibility of
measuring the electron beam polarization using the synchrotron radiation
produced by a magnet karabekov93 . In this paper we examine the feasibility of
a “spin-light” polarimetery technique for measuring longitudinal polarization
of multi-GeV electron beams while building on the 1993 proposal. The simulated
wiggler magnetic field was implemented in a full-fledged Geant4 simulation of
the polarimeter. A polarimeter based on spin-light would provide for a
polarization measurement independent of both Compton and Møller polarimeters.
A relative spin-light polarimeter could also be used in association with
either a Compton or a Møller polarimeter. Highly precise, multiple independent
polarimeters are a must if the ambitious goal of $\sim$ 0.5 % uncertainity in
polarimetry is to be achieved at an Electron Ion Collider (EIC) in order to
meet the experimental demands.
## 2 Spin-Light Charecteristics
The spin-dependent SR distribution as given by Sokolov et. al. [4] , ignoring
higher order effects, is of particular interest as it expresses the
distribution in terms of physical parameters such as spin of the electron -
’$j$’ and the vertical angle - $\psi$ (in the electron’s frame of reference)
[5] . The verticle angle is important as it determines the geometry of the
apparatus besides other design parameters such as position of collimators.
$\displaystyle N_{\gamma}(long)$ $\displaystyle=$
$\displaystyle\frac{9n_{e}}{16\pi^{3}}\frac{e^{2}}{cm_{e}R^{2}}\gamma^{4}\int_{0}^{\infty}\frac{y^{2}dy}{(1+\xi
y)^{4}}\oint d\Omega(1+\alpha^{2})^{2}\times$ (1)
$\displaystyle\left[K^{2}_{2/3}(z)+\frac{\alpha^{2}}{1+\alpha^{2}}K^{2}_{1/3}(z)+j\xi
y\frac{\alpha}{\sqrt{1+\alpha^{2}}}K_{1/3}(z)K_{2/3}(z)\right]$
where $\xi=\frac{3B}{2B_{c}}\gamma$, $B_{c}$ being the magnetic field under
the influenze of which the entire kinetic energy of the electron is expelled
as one SR photon, $y=\frac{\omega_{o}}{\omega_{c}}$, $K_{n}(x)$ are modified
Bessel functions, $n_{e}$ is the number of electrons and,
$z=\frac{\omega}{2\omega_{C}}(1+\alpha^{2})^{3/2}$, and $\alpha=\gamma\psi$.
For an electron that is polarized, the power below (i.e.
$-\frac{\pi}{2}\leq\psi\leq 0$) and above (i.e. $0\leq\psi\leq\frac{\pi}{2}$)
are spin dependent. More importantly the difference between the power radiated
above and power radiated below, called Spin-Light with an assymetry which can
be defined as $A=\frac{\Delta N_{\gamma}}{N_{\gamma}}$, is directly spin
dependent and this opens up the possibility of a direct measurement technique.
Figure 1: Schematic diagram of a differential spin-light polarimeter
## 3 Conceptual Design
The wiggler magnet is at the heart of the setup where the SR photons are
produced and the ionization chambers can be used to characterize the SR in
order to measure the asymmetry. It is important to note the presence of
collimators on the faces of the wiggler magnets in order to prevent
intermixing of the SR light fans. Collimation creates four SR spots, with each
ionization chamber receiving two collimated SR spots. Two major variables in
this setup are the electron beam energy and the wiggler pole strength. In
Figures 2(A) & (B), Spin-Light spectra and the asymmetry are plotted for
various wiggler pole strengths with a $11GeV$ beam and in Figures 2(C) & (D),
the same are plotted for various beam energies with a $4T$ wiggler field. A
numerical integration code is used to generate the SR spectra and asymmetry
using Eq.(1).
Figure 2: (Left- Right): A. Plot of spin light spectra for various pole
strengths from $2T-5T$; B. Plot of asymmetry vs. photon energy for various
pole strengths.; C. Plot of spin light spectra for various beam energies
ranging from $4GeV-12GeV$.; D. Plot of Asymmetry vs. the photon energy for
various beam energies.
### 3.1 Effects of realistic dipole magnetic field with fringes
A field map of the wiggler magnets can be generated by solving Maxwell’s
equations with appropriate boundary conditions. In LANL Poisson SuperFish [18]
, the magnet contours can be easily defined. The field map of the magnet can
then be plotted. Here, the field map at the edge where the electron beam
enters the magnet is presented in Figure 3(A). Note that the beam pipe is
going at the center below the magnet pole. In Figure 3(A), the physical taper
of the cores can be noticed, since it is at the edge of the magnet face. The
singularities seen in Figure 3(A) are the areas where the current cuts the
plane. Also, it is important to note that the entire ‘C’ magnet is not visible
in the field-map, only one half of the ‘C’ magnet is shown in the field map.
The field map obtained can be used in the numerical integration code, in place
of a constant pole strength, to plot the SR spectra and the asymmetry which
are presented in Figures 4(A) & (B). Even though there is a reduction in the
total power output of SR light by introducing a realistic dipole field, the
asymmetry has not changed. This implies that the changes introduces by the
realistic dipoles are negligible.
Figure 3: (Anticlockwise from Top-Left): A. Field map of the dipole face at
the edge of the dipole.; B. Plot of both the $x$ and $y$ components of the
magnetic field on the transverse plane at the the edge of the dipole (Beam
pipe is centered around $15$cm mark along the ’x’ axis).;C. Integrated power
spectra of SR Light at the IC due to Gaussian beam. (The difference between
the profile has been enlarged for clarity); D. Integrated power spectra of SR
Light at the IC due to a point beam.
### 3.2 Effects of Extended Beam Size
Figure 4: (Left- Right): A. Plot showing the SR - Light and Spin - Light power
spectra with a realistic wiggler magnetic field (Power spectra for uniform
magnetic field have also been presented as .; B. Plot of the assymetry with a
realistic wiggler magnetic field.
The effect of having an extended beam size of about $100\mu m$ was studied by
essentially superimposing the SR Spectra generated by each differential
element of the beams weighted with Gaussian distribution in order to make the
extended beam a perfect Gaussian beam. The cumulative spectra for a point beam
was obtained by setting the weighting factor to one. The cumulative spectra
for a Gaussian beam when plotted has approximately the same structure as the
spectra for the point - cross section beam. This is so because the size of the
beam ($R_{beam}=100\mu m$) is small compared to the size of the collimated SR
- Light spot which is about $1mm$ big. For the beam with a point cross
section, the SR - profile is rather ’box’ like at the ionization chamber. When
an extended beam, that is of Gaussian profile, is introduced, the SR - profile
gets a taper which is Gaussian in nature too as illustrated in Figures 3(C) &
(D).
## 4 Spin Light in Geant4 and SR-Spectra
A Geant4 simulation with a geometry as in Figure 1 using EM Extra process list
was constructed. Even though the integrated energy spectrum is validated in
Geant4, it does not contain the angular dependence of SR light. We implement
the angular dependence of SR light at the stacking action level. SR photons
are killed with a probability equal to asymmetry which can be calculated from
standard Geant4 track parameters. This creates an up-down asymmetry in the SR
cones which is vital for this simulation of a Spin-Light polarimeter. As seen
in Figure 5, the simulated SR-Spectra closely matches the physics SR-Spectra
within $1\%$. Spin light component is obtained by subtracting the remaining
tracks with positive momentum (corresponding to $0\leq\psi\leq\frac{\pi}{2}$)
from tracks with negative momentum (corresponding to
$-\frac{\pi}{2}\leq\psi\leq 0$). Figure 5 also shows non-SR events in green
which are significantly small in number compared to the corresponding spin-
light events in a particular photon energy bin.
## 5 Conclusion
Figure 5: Geant4 SR & SL spectra as compared with physics SR & SL spectra.
The figure of merit for such a polarimeter increases with electron beam energy
and the strength of magnetic field used. On the other hand, the SR profile
becomes more compact with increase in electron beam energy. This makes it
difficult for an IC to charecterize the SR profile given that the SR load
increases as fourth power of beam energy. A 3 pole wiggler with a field
strength of $4T$ and a pole length of $10cm$ would be adequate for such a
polarimeter. Locating a reasonable piece of beam-line real estate is however
very challenging. Given that the eRHIC design of EIC involves using a Gatalin
gun [25] at a very high rate, the recovery time of the spin-light IC will
need to be extremely small if every bunch of the electron beam is to be
measured for polarization, a goal which is nearly impossible with this design.
A Spin-light polarimeter is apt for measuring averaged polarization of a
number of beam bunches. A survey of all experiments beings proposed and their
corresponding polarimetry requirements both in terms of precision of
polarimetry required and the rate of measurement will go a long way in helping
pin down the instrument specifications.
## References
* (1) I. P. Karabekov, R. Rossmanith, Proc. of the 1993 PAC, Washington, v. 1, p. 457 (1993); I. P. Karabekov and S. I. Karabekian, Proceedings of 5th European Particle Accelerator Conference (EPAC 96), Sitges, Spain, 10-14 Jun 1996, pp 1743-1745 (1996); A. V. Airapetian, R. O. Avakian, I. P. Karabekov, E. L. Saldin, and M. V. Yurkov, Proc. of the SPIN-96, Amsterdam, The Netherlands, Vol.1, p762 (1996).
* (2) A. A. Sokolov, N. P. Klepikov and I. M. Ternov, JETF 23, 632 (1952).
* (3) A. A. Sokolov and I. .M. Ternov, Radiation from Relativistic Electrons, A.I.P. Translation Series, New York (1986) ; I. M. Ternov, Physics - Uspekhi 38, 409 (1995).
* (4) K. Sato, J. of Synchrotron Rad., 8, 378 (2001).
* (5) B. Norum, CEBAF Technical note, TN-0019 (1985).
* (6) M. Sands, SLAC Technical note, SLAC-121 (1970).
* (7) Poisson SuperFish 2D EM Solver, laacg1.lanl.gov/laacg/services/sfu_04_04_03.phtml, 2007.
* (8) V. N. Litvinenko, Gatling Gun: High Average Polarized Current Injector for eRHIC, EIC BNL Whitepapers (2012).
|
arxiv-papers
| 2013-09-11T22:17:08 |
2024-09-04T02:49:50.795819
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Prajwal Mohanmurthy and Dipangkar Dutta",
"submitter": "Prajwal Mohanmurthy",
"url": "https://arxiv.org/abs/1309.2988"
}
|
1309.3065
|
UT–13–34
Probing Bino Contribution to Muon $g-2$
Motoi Endo, Koichi Hamaguchi,
Teppei Kitahara, and Takahiro Yoshinaga***Research Fellow of the Japan Society
for the Promotion of Science
Department of Physics, University of Tokyo, Tokyo 113-0033, Japan
We study SUSY models in which Bino contributions solve the muon $g-2$ anomaly.
The contributions are enhanced by large left-right mixing of the smuons.
However, it is constrained by the vacuum stability condition of the
slepton–Higgs potential. Therefore, there are upper bounds on masses of
sleptons and Bino. When the slepton soft masses are universal, the upper bound
on the smuon mass becomes $330~{}(460)\,{\rm GeV}$ in order to solve the $g-2$
anomaly at the $1\sigma~{}(2\sigma)$ level. It is within the reach of LHC and
ILC. If the stau is heavier than the smuon, the bound can be as large as
$1.4~{}(1.9)\,{\rm TeV}$. Such non-universal slepton mass spectrum generically
predicts too large LFV/CPV. We show that the models are expected to be probed
by LHC/ILC and LFV/CPV complementarily in future.
## 1 Introduction
The anomalous magnetic moment of the muon, $a_{\mu}=(g-2)_{\mu}/2$, (muon
$g-2$) has been measured very precisely [1]. It is compared with the Standard
Model (SM) prediction, and the latest result is
$\Delta a_{\mu}\equiv a_{\mu}({\rm exp})-a_{\mu}({\rm SM})=(26.1\pm 8.0)\times
10^{-10},$ (1)
where Ref. [2] is referred to for contributions of the hadronic vacuum
polarization, and the hadronic light-by-light contribution is from Ref. [3].
Similarly, Ref. [4] provides $\Delta a_{\mu}=(28.6\pm 8.0)\times 10^{-10}$.
Therefore, there is discrepancy at more than $3\sigma$ confidence levels. It
is noticed that the difference is as large as SM electroweak contributions,
$a_{\mu}({\rm EW})=(15.36\pm 0.1)\times 10^{-10}$ [5, 6]. If this is a
signature of physics beyond the SM, and if the new physics exists in the TeV
scale, new physics contributions to the muon $g-2$ are necessarily enhanced by
some mechanisms, because they are naively estimated as $a_{\mu}({\rm
NP})\sim(\alpha_{\rm NP}/4\pi)\times(m_{\mu}^{2}/m_{\rm NP}^{2})$, which is
required to be comparable to $a_{\mu}({\rm EW})$.
The supersymmetry (SUSY) is a good candidate for such new physics models. The
model can provide sizable contributions to the muon $g-2$ [7, 8, 9], which are
enhanced by $\tan\beta=\langle H_{u}\rangle/\langle H_{d}\rangle$. The muon
$g-2$ anomaly is solved if the superparticles (muonic sleptons, neutralinos
and/or charginos) have a mass around $\mathcal{O}(100)\,{\rm GeV}$ for
$\tan\beta=\mathcal{O}(10)$. They are light enough to be produced at the LHC
experiments. Recently, LHC phenomenology of the superparticles that are
relevant for the muon $g-2$ has been studied in Ref. [10]. It has been shown
that they can be discovered in the near future in most of the parameter
regions, especially in regions where the SUSY contributions to the muon $g-2$
are dominated by chargino–sneutrino diagrams. However, the searches rely on
the assumption that the Wino is light. This assumption is not necessary to
explain the muon $g-2$ discrepancy, when the SUSY contribution is mainly from
Bino–smuon diagrams. If the Wino is heavy, collider searches differ
significantly from those in Ref. [10]. In this paper, we study searches for
the models in which only the Bino and the left- and right-handed smuons are
light, while the other superparticles including the Wino are decoupled from
the LHC sensitivity.
The Bino–smuon contribution is enhanced when the left-right mixing of the
smuon is large. If the left-right mixing were allowed to be arbitrarily large,
the superparticles could be extremely heavy while keeping the contribution to
the muon $g-2$, thereby escaping any collider searches. However, too large
left-right mixing spoils stability of the electroweak vacuum. Hence, the
superparticle masses are bounded from above. We will show that, when staus and
smuons have comparable masses to each others, upper bounds on the slepton
masses are within the reach of LHC/ILC. If the staus are much heavier than
smuons, the mass bounds are relaxed, while lepton flavor violations (LFV) and
CP violations (CPV) generically become too large. Therefore, we will show that
almost all the parameter regions can be tested in future complementarily by
LHC/ILC and LFV/CPV, if the muon $g-2$ anomaly is solved by the Bino and the
left- and right-handed smuons.
The rest of the paper is organized as follows. The mass spectrum is provided
in Sec. 2. The SUSY contributions to the muon $g-2$ and the vacuum stability
conditions are explained in Sec. 3. In Sec. 4, experimental searches are
studied. Sec. 5 is devoted to the conclusion.
## 2 Mass spectrum
We assume that the SUSY contributions to the muon $g-2$ are dominated by the
Bino–smuon contribution. The Bino and the left- and right-handed smuons
contribute to the diagram. It is enhanced by the left-right mixing of the
smuon, which is determined by the muon Yukawa coupling constant, the Higgsino
mass parameter, $\mu$, and $\tan\beta$. Details will be discussed in the next
section. Since too large $\tan\beta$ spoils perturbativity of the down-type
Yukawa interactions, $\mu$ is favored to be large to explain the muon $g-2$
discrepancy (1). Thus, we focus on large Higgsino mass regions. The Winos are
also supposed to be decoupled.
The collider searches depend on the slepton mass spectrum. The left- and
right-handed selectron masses are assumed to be degenerate with those of the
smuons, respectively. The following conclusion does not depend on this
assumption.111 If the selectron masses are non-universal in the first two
generations, the LFV bound is severe. See the discussions in Sec. 4.2. On the
other hand, the vacuum stability condition and the LFV/CPV bounds depend on
the stau masses (see Sec. 3.2 and 4.2).
All colored superparticles are set to be very heavy. In fact, none of them
have been discovered at LHC. The Higgs boson mass of $126\,{\rm GeV}$ favors
the scalar top masses to be $\mathop{\mathcal{O}}(10-100)\,{\rm TeV}$, if the
trilinear coupling of the top squark is not large. Similarly, the heavy Higgs
bosons of the two Higgs doublets are assumed to be heavy. In this paper, all
of them are considered to be decoupled.
Consequently, we consider the low-energy effective theory, in which only the
following superparticles are light,
$\tilde{B},~{}\tilde{\ell}_{L},~{}\tilde{\ell}_{R}.$ (2)
Here, $\tilde{\ell}$ denotes the selectron and the smuon (and the stau,
depending on the mass spectrum). In addition, the Higgsinos can contribute to
the effective Lagrangian, if they are light due to the vacuum stability
condition. The model parameters are
$M_{1},m_{\tilde{\ell}_{L}}^{2},m_{\tilde{\ell}_{R}}^{2},m_{\tilde{\ell}_{LR}}^{2}.$
(3)
Here, $M_{1}$ is the Bino mass. Since the Higgsinos (and the Wino) are heavy,
the lightest neutralino is almost composed of the Bino. On the other hand,
$m_{\tilde{\ell}_{L}}^{2}$ and $m_{\tilde{\ell}_{R}}^{2}$ are soft SUSY-
breaking masses of the left- and right-handed sleptons, respectively.
$m_{\tilde{\ell}_{LR}}^{2}$ is off-diagonal components of the slepton mass
matrices. This controls the observables and constraints that are studied in
this paper.
Before proceeding to the phenomenology, let us comment on the left-right
mixing of the smuon. It includes the scalar trilinear coupling of the muon,
$A_{\mu}$, as well as $\mu$ and $\tan\beta$. The Bino–smuon contribution could
be enhanced by $A_{\mu}$ with $\mu\tan\beta$ kept small. However, this
requires $A_{\mu}$ to be extraordinary large. If the trilinear coupling is
universal among the matter scalar fermions, this implies extremely large
trilinear coupling for the stop sector, resulting in either too large Higgs
boson mass or rapid decays of our electroweak vacuum into charge/color
breaking vacua. In this paper, $A_{\mu}$ is set to be zero for simplicity, and
the left-right mixing of the smuon is determined by $\mu$ and $\tan\beta$.
## 3 Muon $g-2$ and Vacuum Stability
### 3.1 Muon $g-2$
Figure 1: The Bino–smuon contribution to the magnetic dipole operators.
The Bino–smuon contribution to the muon $g-2$ is represented by Fig. 1, where
the lepton $\ell$ is the muon. In the mass insertion approximation, it is
estimated as (cf. [9])
$\displaystyle a_{\mu}({\rm SUSY})$ $\displaystyle=-(1+\delta^{\rm
2loop})\frac{\alpha_{Y}}{4\pi}\frac{m_{\mu}M_{1}m_{\tilde{\mu}_{LR}}^{2}}{m^{2}_{\tilde{\mu}_{L}}m^{2}_{\tilde{\mu}_{R}}}\,f_{N}\left(\frac{m^{2}_{\tilde{\mu}_{L}}}{M^{2}_{1}},\frac{m^{2}_{\tilde{\mu}_{R}}}{M^{2}_{1}}\right)$
$\displaystyle=\frac{1+\delta^{\rm
2loop}}{1+\Delta_{\mu}}\frac{\alpha_{Y}}{4\pi}\frac{m_{\mu}^{2}M_{1}\mu}{m^{2}_{\tilde{\mu}_{L}}m^{2}_{\tilde{\mu}_{R}}}\tan\beta\cdot
f_{N}\left(\frac{m^{2}_{\tilde{\mu}_{L}}}{M^{2}_{1}},\frac{m^{2}_{\tilde{\mu}_{R}}}{M^{2}_{1}}\right)$
$\displaystyle\simeq 1.5\times 10^{-9}~{}\frac{1+\delta^{\rm
2loop}}{1+\Delta_{\mu}}\left(\frac{\tan\beta}{10}\right)\left(\frac{(100\,{\rm
GeV})^{2}}{m_{\tilde{\mu}_{L}}^{2}m_{\tilde{\mu}_{R}}^{2}/M_{1}\mu}\right)\left(\frac{f_{N}}{1/6}\right),$
(4)
at the leading order of $\tan\beta$. The loop function $f_{N}(x,y)$ is defined
as
$f_{N}(x,y)=xy\left[\frac{-3+x+y+xy}{(x-1)^{2}(y-1)^{2}}+\frac{2x\log
x}{(x-y)(x-1)^{3}}-\frac{2y\log y}{(x-y)(y-1)^{3}}\right],$ (5)
which satisfies $f_{N}(1,1)=1/6$ and $0\leq f_{N}(1,1)\leq 1$. It is noticed
that the contribution is proportional to the left-right mixing,
$\displaystyle
m_{\tilde{\ell}_{LR}}^{2}=-Y_{\ell}\,v_{u}\mu=-\frac{m_{\ell}}{1+\Delta_{\ell}}\mu\tan\beta,$
(6)
where $v_{u}$ is the vacuum expectation value (VEV) of the up-type Higgs
field. Note that the heavy Higgs bosons are decoupled. In the last equation,
$\Delta_{\ell}$ is a correction to the lepton Yukawa coupling constant [11].
It appears when the lepton Yukawa coupling is matched to the physical lepton
mass, $m_{\ell}$, or the Yukawa coupling in SM. In the low-energy effective
theory, it becomes
$\displaystyle\Delta_{\ell}=\frac{\alpha_{Y}}{4\pi}M_{1}\mu\tan\beta\cdot
I(M_{1}^{2},m_{\tilde{\ell}_{L}}^{2},m_{\tilde{\ell}_{R}}^{2}),$ (7)
where Higgsino diagrams are discarded, because they are suppressed in large
Higgsino mass regions. Also, the terms that are not enhanced by $\tan\beta$
are omitted. The loop function is defined as
$\displaystyle I(a,b,c)=-\frac{ab\ln a/b+bc\ln b/c+ca\ln
c/a}{(a-b)(b-c)(c-a)},$ (8)
which satisfies $f(a,a,a)=1/2a$. In particular, when $\mu\tan\beta$ is very
large, $\Delta_{\mu}$ becomes as large as or larger than
$\mathop{\mathcal{O}}(0.1-1)$.
The correction $\delta^{\rm 2loop}$ denotes leading contributions of two loop
diagrams. It is estimated as
$\displaystyle 1+\delta^{\rm 2loop}=\left(1-\frac{4\alpha}{\pi}\ln\frac{m_{\rm
soft}}{m_{\mu}}\right)\left[1+\frac{1}{4\pi}\left(2\alpha_{Y}\Delta
b+\frac{9}{4}\alpha_{2}\right)\ln\frac{M_{\rm soft}}{m_{\rm soft}}\right].$
(9)
In the right-hand side, the first bracket is QED corrections to the muon $g-2$
[12], or renormalization group contributions to the effective operator of the
magnetic dipole operator from the (smuon or Bino) soft mass scale, $m_{\rm
soft}$, to the muon mass scale. In the numerical analysis, $m_{\rm soft}$ is
chosen to be the smuon mass. This correction is $\sim 10\%$. Non-logarithmic
terms evaluated in Ref. [12] are found to be very small in the parameter
regions of this paper.
The second bracket in Eq.(9) is corrections to the Bino couplings with the
smuons. When SUSY is exact, the gaugino coupling is equal to the gauge
coupling constant. This equality is violated after heavy superparticles are
decoupled at a scale, $M_{\rm soft}$. The Bino–muon–smuon interactions are
$\displaystyle\mathcal{L}_{\text{int}}=-\frac{1}{\sqrt{2}}\tilde{g}_{L}\,\overline{\tilde{B}}\mu_{L}\,\tilde{\mu}^{*}_{L}+\sqrt{2}\tilde{g}_{R}\,\overline{\tilde{B}}\mu_{R}\,\tilde{\mu}_{R}^{*}+{\rm
h.c.},$ (10)
where the coefficients are
$\displaystyle\tilde{g}_{L}$ $\displaystyle=g_{Y}+\delta\tilde{g}_{L}\simeq
g_{Y}\left[1+\frac{1}{4\pi}\left(\alpha_{Y}\Delta
b+\frac{9}{4}\alpha_{2}\right)\ln\frac{M_{\rm soft}}{m_{\rm soft}}\right],$
(11) $\displaystyle\tilde{g}_{R}$
$\displaystyle=g_{Y}+\delta\tilde{g}_{R}\simeq
g_{Y}\left[1+\frac{\alpha_{Y}}{4\pi}\Delta b\,\ln\frac{M_{\rm soft}}{m_{\rm
soft}}\right].$ (12)
Here, the terms of $\alpha_{Y}$ are corrections to the Bino self-energy, which
are SUSY analog of the oblique corrections [13, 14, 15]. The U(1)Y gauge
coupling constant, $g_{Y}$, is evaluated at $m_{\rm soft}$, and the
corrections are represented by the difference between the beta functions of
the U(1)Y gauge and Bino couplings. In our setup, the SM particles and the
sleptons contribute to the beta functions. Thus, the coefficient becomes
$\Delta b=41/6-n_{\rm slepton}$, where $n_{\rm slepton}$ is number of the
generations of light sleptons; for instance, $n_{\rm slepton}=3$ if all the
sleptons are light, and $n_{\rm slepton}=2$ when the staus are decoupled. On
the other hand, the term of $\alpha_{2}$ is non-oblique corrections after the
Wino decoupled (cf. Ref. [16]). The weak boson couples only to the left-handed
(s)leptons at the Bino–muon–smuon vertices. Since the Bino–smuon contribution
to the muon $g-2$ includes both $\tilde{g}_{L}$ and $\tilde{g}_{R}$, the loop
correction (9) is obtained. It yields a correction of $5-10\%$ for $M_{\rm
soft}=10-100\,{\rm TeV}$ with $m_{\rm soft}\sim 100\,{\rm GeV}$. In the
analysis, only the logarithmic contributions are included. Non-logarithmic
terms are expected to be suppressed.
Figure 2: The SUSY contributions to the muon $g-2$ and the vacuum stability
bounds are shown. In the orange (yellow) regions, the muon $g-2$ discrepancy
(1) is explained at the $1\sigma$ $(2\sigma)$ level. The black solid lines
represent the upper bound on $\mu$ from the vacuum stability bound of the
stau-Higgs potential. The stau masses are set to be
$m_{\tilde{\tau}}/m_{\tilde{\mu}}=1,2,5,10$ and 20 from bottom to top. Above
the black dashed line, $\Delta_{\mu}$ becomes larger than unity. The
parameters are $M_{1}=m_{\tilde{\mu}_{L}}=m_{\tilde{\mu}_{R}}$, $\tan\beta=40$
and $M_{\rm soft}=30\,{\rm TeV}$. In the right panel, a part of the parameter
region of the left panel is magnified.
In Eq. (9), the corrections that can be as large as or larger than
$\mathop{\mathcal{O}}(10)\%$ are included. Other two loop contributions are
unknown or expected to be small. In Refs. [17, 18], SUSY corrections to the SM
one-loop diagrams are evaluated. When the superparticles except those of (2)
are decoupled, the contributions in Refs. [17, 18] are negligibly small. Other
corrections that have not yet been estimated include electroweak and SUSY two-
loop contributions to SUSY one-loop diagrams. They might provide $\delta
a_{\mu}\sim 10^{-10}$, according to Ref. [19]. In addition, non-logarithmic
corrections to $\delta^{\rm 2loop}$ could be a few percents of the SUSY one-
loop contributions, similarly to the discussions in Ref. [19].
Apart from the Bino–smuon diagrams, there are other one-loop contributions to
the muon $g-2$. The Bino–Higgsino–smuon contribution can be
$\lesssim\mathop{\mathcal{O}}(10^{-10})$ when the Higgsinos are light due to
the vacuum stability bound (Sec. 3.2). It is included in the numerical
analysis for completeness.222 This contribution can dominate the SUSY
contributions to the muon $g-2$, when $\mu$ is small while decoupling the
Wino. Since they are enhanced only by $\tan\beta$, superparticles are required
to be light to explain (1). They are detectable in colliders. In particular,
the Higgsino production can be significant. On the other hand, the
chargino–muon sneutrino contributions are less than
$\mathop{\mathcal{O}}(10^{-11})$ for $M_{2}>10\,{\rm TeV}$, i.e., negligible.
In Fig. 2, contours of the SUSY contributions to the muon $g-2$ are shown. The
horizontal and vertical axises are the lightest smuon mass,
$m_{\tilde{\mu}_{1}}$, and $\mu$, respectively. The parameters are set as
$M_{1}=m_{\tilde{\mu}_{L}}=m_{\tilde{\mu}_{R}}$, $\tan\beta=40$ and $M_{\rm
soft}=30\,{\rm TeV}$. In the orange (yellow) regions, the SUSY contributions
explain the muon $g-2$ discrepancy (1) at the $1\sigma$ $(2\sigma)$ level. It
is found that they are enhanced by large $\mu$, and the smuon masses can be
$1\,{\rm TeV}$ for $\mu=\mathop{\mathcal{O}}(10-100)\,{\rm TeV}$. This is
contrasted to the chargino–muon sneutrino contributions to the muon $g-2$,
where $\mu$ is favored to be small [10]. On the other hand, detailed
dependences on the superparticle mass spectrum are determined by the loop
function (5) and the vacuum stability condition. They will be discussed in the
next subsection.
### 3.2 Vacuum stability
As shown in Sec. 3.1, the Bino–smuon contribution to the muon $g-2$ is
enhanced by a large left-right mixing of the smuon. However, too large mixing
spoils the stability of the electroweak vacuum. The trilinear coupling of the
sleptons and the SM-like Higgs boson is given by
$\displaystyle V$
$\displaystyle\simeq\frac{1}{\sqrt{2}v}m_{\tilde{\ell}_{LR}}^{2}\tilde{\ell}_{L}^{*}\tilde{\ell}_{R}h^{0}+{\rm
h.c.}$
$\displaystyle=-\frac{m_{\ell}}{\sqrt{2}v(1+\Delta_{\ell})}\mu\tan\beta\cdot\tilde{\ell}_{L}^{*}\tilde{\ell}_{R}h^{0}+{\rm
h.c.},$ (13)
where $v\simeq 174\,{\rm GeV}$ is the Higgs VEV. As the trilinear coupling
increases, disastrous charge-breaking minima in the scalar potential become
deeper, and our electroweak vacuum could decay to them. By requiring that the
lifetime of the electroweak vacuum should be longer than the age of the
Universe, $m_{\tilde{\ell}_{LR}}^{2}$ is constrained.
The vacuum stability conditions have been studied. The fitting formula of the
stability condition is obtained as
$\displaystyle\left|m_{\tilde{\ell}_{LR}}^{2}\right|$
$\displaystyle\leq\eta_{\ell}\bigg{[}1.01\times 10^{2}\,{\rm
GeV}\sqrt{m_{\tilde{\ell}_{L}}m_{\tilde{\ell}_{R}}}+1.01\times 10^{2}\,{\rm
GeV}(m_{\tilde{\ell}_{L}}+1.03m_{\tilde{\ell}_{R}})$
$\displaystyle-2.27\times$ $\displaystyle 10^{4}\,{\rm
GeV}^{2}+\frac{2.97\times 10^{6}\,{\rm
GeV}^{3}}{m_{\tilde{\ell}_{L}}+m_{\tilde{\ell}_{R}}}-1.14\times 10^{8}\,{\rm
GeV}^{4}\left(\frac{1}{m^{2}_{\tilde{\ell}_{L}}}+\frac{0.983}{m^{2}_{\tilde{\ell}_{R}}}\right)\bigg{]}.$
(14)
This is consistent with the result of Ref. [20] in a case of the stau–Higgs
potential. The vacuum decay rate is evaluated by the bounce method [21]. The
public package CosmoTransitions 1.0.2 [23] is used for the numerical
evaluation. The scalar potential of the left- and right-handed sleptons and
the SM-like Higgs boson is analyzed at the zero temperature.333 Thermal
corrections can change the transition rate. They could be significant when the
staus are light, for instance, for $m_{\tilde{\tau}_{1}}\lesssim 200\,{\rm
GeV}$ [22]. The Higgs potential is set to reproduce the mass of $126\,{\rm
GeV}$. Only the renormalizable terms are taken into account. Higher-
dimensional terms depend on the superparticles that decouple at $M_{\rm
soft}$. Since $M_{\rm soft}$ is very large, their contributions to (14) are
considered to be small. The fitting formula reproduces results of
CosmoTransitions at better than the 1% level.
The fitting formula (14) is universal for all the slepton flavors up to
corrections, $\eta_{\ell}\sim 1$. Given the soft scalar masses, the bound is
independent of the flavor except through Yukawa interactions in the quartic
terms of the scalar potential. Since the quartic terms are dominated by gauge
interactions, $\eta_{\ell}$ changes little around unity. It also depends on
$\tan\beta$ when the Yukawa coupling is large. In practice, $\eta_{\tau}=1$ is
set for the stau with $\tan\beta/(1+\Delta_{\tau})=70$. Numerical result of
$\eta_{\tau}$ for the stau with various $\tan\beta$ is found in Fig. 2 of Ref.
[20]. For instance, it is $\eta_{\tau}\simeq 0.94$ for
$\tan\beta/(1+\Delta_{\tau})=40$. On the other hand, we obtain
$\eta_{\ell}\simeq 0.88$ for the smuons and the selectrons, which is
independent of the flavor and $\tan\beta$ because of small Yukawa couplings.
The most severe constraint on $\mu\tan\beta$ is obtained from the stability
condition of the stau-Higgs potential. This is because the left-right mixing
is proportional to the Yukawa coupling. In Fig. 2, the upper limits are shown
for $\tan\beta=40$. Combined with the muon $g-2$, the smuon masses are bounded
from above for given stau masses. For
$m_{\tilde{e}}=m_{\tilde{\mu}}=m_{\tilde{\tau}}$, the lightest smuon is
limited to be $m_{\tilde{\mu}_{1}}\lesssim 300~{}(420)\,{\rm GeV}$ at the
$1\sigma$ ($2\sigma$) level of the muon $g-2$.
Figure 3: The SUSY contributions to the muon $g-2$ as a function of the
lightest smuon mass, $m_{\tilde{\mu}_{1}}$, and the lightest neutralino mass,
$m_{\tilde{\chi}^{0}_{1}}$. In the orange (yellow) regions, the muon $g-2$
discrepancy (1) is explained at the $1\sigma$ $(2\sigma)$ level. The left-
right mixing is maximized under the vacuum stability condition. The parameters
are $m_{\tilde{\ell}_{L}}=m_{\tilde{\ell}_{R}}$, $\tan\beta=40$ and $M_{\rm
soft}=10\,{\rm TeV}$. The stau soft masses are equal to those of the
selectrons and smuons. The region below the green line is excluded by LHC.
The SUSY contributions to the muon $g-2$ depend on superparticle mass spectra
mainly through the loop function (5). In Fig. 3, contours of the muon $g-2$
are displayed for various smuon and neutralino masses. In the orange (yellow)
regions, the muon $g-2$ discrepancy is explained at the $1\sigma$ ($2\sigma$)
level. The left-right mixing is maximized with satisfying the vacuum stability
condition. Here, the stau soft mass is supposed to be degenerate with that of
the smuon. It is found that the lightest smuon mass can be $330~{}(460)\,{\rm
GeV}$ at the $1\sigma$ ($2\sigma$) level of the muon $g-2$, when the lightest
neutralino mass, $m_{\tilde{\chi}^{0}_{1}}$, is almost a half of
$m_{\tilde{\mu}_{1}}$. Note that $m_{\tilde{\chi}^{0}_{1}}$ is almost equal to
$M_{1}$. On the other hand, $m_{\tilde{\chi}^{0}_{1}}\simeq 560\,{\rm GeV}$ is
realized for $m_{\tilde{\mu}_{1}}\simeq 250\,{\rm GeV}$ at the $2\sigma$ level
of the muon $g-2$.
For a fixed value of the lightest smuon mass, the SUSY contributions to the
muon $g-2$ is maximized when $m_{\tilde{\mu}_{L}}$ is equal to
$m_{\tilde{\mu}_{R}}$. Eq. (4) depends on the left- and right-handed smuon
masses as $f_{N}(x,y)/xy$ with $x=m_{\tilde{\mu}_{L}}^{2}/M_{1}^{2}$ and
$y=m_{\tilde{\mu}_{R}}^{2}/M_{1}^{2}$. When $xy$ is fixed, it is maximized for
$x=y$ and rapidly decreases for $x\neq y$. On the other hand, the vacuum
stability condition (14) is relaxed for $m_{\tilde{\ell}_{L}}\neq
m_{\tilde{\ell}_{R}}$. Since the former dependence is stronger than the
latter, the SUSY contributions to the muon $g-2$ become largest when
$m_{\tilde{\mu}_{L}}$ is equal to $m_{\tilde{\mu}_{R}}$ for given $M_{1}$ and
$\tan\beta$. In other words, the orange/yellow regions in Fig. 3, i.e., the
parameter regions favored by the muon $g-2$, shrink toward smaller
superparticle masses for $m_{\tilde{\mu}_{L}}\neq m_{\tilde{\mu}_{R}}$.
Figure 4: The left-right mixing of the smuon (blue solid) and the stau (red
solid) as a function of $\mu$. The blue (red) dashed line represents the
vacuum meta-stability bound with $\eta=0.88~{}(0.94)$. Above the black dashed
line, sleptons become tachyonic. Below the blue dotted line, $\Delta_{\mu}<1$
is satisfied. The parameters are
$m_{\tilde{\ell}_{L}}=m_{\tilde{\ell}_{R}}=2M_{1}=500\,{\rm GeV}$,
$\tan\beta=40$ and $M_{\rm soft}=30\,{\rm TeV}$.
The vacuum stability bound from the stau-Higgs potential is relaxed for
heavier staus. When they are decoupled, the constraint almost disappears, and
larger $\mu$ is allowed. As $\mu$ increases, the left-right mixing (6)
approaches to a maximal value. It is determined by the soft scalar masses as
$\displaystyle
m_{\tilde{\ell}_{LR}}^{2}=-\frac{m_{\ell}}{1+\Delta_{\ell}}\mu\tan\beta~{}\to~{}-\frac{4\pi}{\alpha_{Y}}\frac{m_{\ell}}{M_{1}}\frac{1}{I(M_{1}^{2},m_{\tilde{\ell}_{L}}^{2},m_{\tilde{\ell}_{R}}^{2})}~{}~{}~{}(\text{for}~{}\mu\to\infty).$
(15)
Note that the correction $\Delta_{\ell}$ is much larger than unity in this
limit. In Fig. 4, the left-right mixings of the smuon and the stau,
$|m_{\tilde{\ell}_{LR}}^{2}|$, are plotted as a function of $\mu$, by the blue
and red solid lines, respectively. Since $m_{\tilde{\ell}_{LR}}^{2}$ is
proportional to the Yukawa coupling, $|m_{\tilde{\tau}_{LR}}^{2}|$ is larger
than $|m_{\tilde{\mu}_{LR}}^{2}|$. When the staus are light, $\mu$ is bounded
from above by the stability condition of the stau–Higgs potential (red dashed
line). For heavier staus, the red dashed line (as well as the stau tachyonic
bound or the black dashed line) is lifted up, and larger $\mu$ is allowed.
Even if the staus are decoupled, the left-right mixing is limited by the
vacuum stability condition of the smuon-Higgs potential. When $M_{1}$ is
smaller than the smuon mass with $m_{\tilde{\mu}_{L}}=m_{\tilde{\mu}_{R}}$,
the vacuum constraint (14) is more severe than Eq. (15). It can be estimated
that the lightest smuon can be as large as $1.4~{}(1.9)\,{\rm TeV}$ at the
$1\sigma~{}(2\sigma)$ level of the muon $g-2$.444 In this limit, the situation
is essentially the same as the soft Yukawa coupling models [24, 25]. In this
paper, the vacuum stability condition is analyzed seriously and is found to
restrict the SUSY contributions to the muon $g-2$, resulting in tighter bound
on the smuon mass than Ref. [25]. Since it is very weak, the bound does not
appear in the parameter range of Fig. 2.
It is noticed that $\Delta_{\ell}$ is larger than unity when $\mu$ is very
large. Then, the radiative correction exceeds the tree level contribution to
the lepton mass. For $\Delta_{\ell}\gg 1$, the Yukawa coupling constant is
suppressed by $\Delta_{\ell}$, as observed by Eq. (6). This may be disfavored
by a naturalness argument [26]. In Fig. 2, $\Delta_{\mu}$ becomes larger than
unity above the black dashed line, and above the blue dotted line in Fig. 4.
## 4 Searches
Mass spectrum | Smuon | Vacuum | LHC/ILC | LFV/EDM
---|---|---|---|---
$m_{\tilde{e}}=m_{\tilde{\mu}}=m_{\tilde{\tau}}$ (Sec. 4.1) | $<330/460\,{\rm GeV}$ | $\tilde{\tau}$ | ✓ |
$m_{\tilde{e}}=m_{\tilde{\mu}}<m_{\tilde{\tau}}$ (Sec. 4.2) | $<1.4/1.9\,{\rm TeV}$ | $\tilde{\tau}$ or $\tilde{\mu}$ | | ✓
Table 1: Summary of searches. The lightest smuon mass is restricted to explain
the muon $g-2$ discrepancy at the $1\sigma/2\sigma$ level. The left-right
mixing is limited by the vacuum stability condition of the stau– or
smuon–Higgs potential, depending on the stau masses. Models of the universal
mass spectrum can be tested by LHC/ILC, while those of the non-universal
spectrum predict large LFV/EDM.
In the previous section, it has been shown that when the muon $g-2$ anomaly
(1) is solved by the Bino–smuon contribution (4), soft masses of the Bino and
the smuon are bounded from above by the vacuum stability condition (14). The
result is summarized in Tab. 1. In this section, we study experimental status
and future prospects to search for such SUSY models. The mass bounds depend on
the slepton mass spectrum. When the stau is degenerate with the smuon, a tight
constraint is imposed by the stau stability condition. In Sec. 4.1, it will be
discussed that the limit is strong enough for the superparticles to be
detectable directly in colliders. On the other hand, if the staus are heavier
or decoupled, the stability bound is relaxed. The superparticle masses can
exceed the collider sensitivities. In Sec. 4.2, it will be argued that such
hierarchical mass spectrum can be probed by LFV and CPV.
### 4.1 Universal slepton mass
In this section, we study collider searches for the SUSY models with the
universal slepton mass spectrum. Here, the left- (right-) handed selectron,
smuon and stau have a common soft SUSY-breaking mass,
$\displaystyle
m_{\tilde{e}_{L}}=m_{\tilde{\mu}_{L}}=m_{\tilde{\tau}_{L}},~{}~{}~{}m_{\tilde{e}_{R}}=m_{\tilde{\mu}_{R}}=m_{\tilde{\tau}_{R}}.$
(16)
Then, the vacuum stability condition from the stau–Higgs potential restricts
slepton masses up to $330-460\,{\rm GeV}$ to solve the muon $g-2$ anomaly.
This is within the reach of the LHC or ILC sensitivity. In fact, some of the
parameter regions have already been excluded by LHC, as shown later.
Collider signatures depend on the lightest superparticle (LSP). In the case of
universal soft slepton masses, either the lightest neutralino or the lightest
stau is LSP among the MSSM particles.555 The lightest stau is lighter than
sneutrinos when $\mu$ is large. In the latter case, the stau is likely to be
long-lived.666 The stau could decay in detectors, for instance, through
R-parity violations. The SUSY signatures depend on decay channels. Such cases
are not discussed here. The (meta-) stable staus leave charged tracks in
detectors. Such signatures have been studied by ATLAS [27] and CMS [28]. The
CMS constraint on the cross section of the stau direct production provides the
95% CL exclusion limit, $m_{\tilde{\tau}_{1}}>339\,{\rm GeV}$. If this is
imposed in addition to the vacuum stability condition, all the parameter
regions of the long-lived stau are excluded in Fig. 3.
Figure 5: Same as Fig. 3, but with a condition that the neutralino is LSP,
$m_{\tilde{\tau}_{1}}>m_{\tilde{\chi}^{0}_{1}}$.
In Fig. 5, we show the SUSY contributions to the muon $g-2$ when the left-
right mixing is maximized under the two conditions; (i) the vacuum stability
and (ii) the neutralino LSP, $m_{\tilde{\tau}_{1}}>m_{\tilde{\chi}^{0}_{1}}$.
Then, SUSY signature is (opposite-sign same-flavor) di-lepton with large
missing transverse energy. Selectrons and smuons are produced by collisions.
They subsequently decay into the lightest neutralino and a partner lepton.
Recently, this signature was studied by ATLAS [29] and CMS [30]. In
particular, the 95% CL exclusion limit is obtained for
$m_{\tilde{\mu}_{L}}=m_{\tilde{\mu}_{R}}$ by ATLAS. The result is shown by the
green solid line in Fig. 5 (and in Fig. 3). Here, the region below the line is
excluded. In detail, the left-right mixing is negligible in the ATLAS
analysis. However, the left- and right-handed smuons maximally mix with each
other in Fig. 5. The total cross section of the smuon productions decreases by
10% compared to the ATLAS setup. Since it is sufficiently small, the exclusion
limit is considered to be almost the same as the ATLAS result. As a result, it
is found that a large fraction of the $1\sigma$ parameter region of the muon
$g-2$ is already excluded.777 In the ATLAS analysis,
$m_{\tilde{e}}=m_{\tilde{\mu}}$ is assumed. The selectron production provides
almost the same constraint as the smuon [29]. Thus, if selectrons are
decoupled, the mass bound becomes weaker. However, LFV/CPV constraints are
severe, as will be mentioned in Sec. 4.2. The sensitivity will be improved by
the upgrade of the energy and the luminosity. In Fig. 6, the total cross
section of the smuon productions are presented for $\sqrt{s}=8\,{\rm TeV}$ and
$14\,{\rm TeV}$. The cross section is estimated at the leading order. For
instance, it becomes $1\,{\rm fb}$ for the lightest smuon mass of $330\,{\rm
GeV}$ at $\sqrt{s}=8\,{\rm TeV}$, which corresponds to the current LHC bound.
The same cross section is obtained for $450\,{\rm GeV}$ at $14\,{\rm TeV}$.
Studies for the future sensitivity are required.
Figure 6: The total cross section of the smuon productions at LHC with
$\sqrt{s}=8\,{\rm TeV}$ (dashed) and $14\,{\rm TeV}$ (solid). The parameters
satisfy $m_{\tilde{\ell}_{L}}=m_{\tilde{\ell}_{R}}=2M_{1}$. The left-right
mixing is maximized under the vacuum stability condition with $\tan\beta=40$.
When the neutralino is LSP, the stau productions have also been studied in LEP
and LHC. If their mass difference is larger than $15\,{\rm GeV}$, the stau
mass is constrained to be larger than $81.9\,{\rm GeV}$ by LEP [31]. LHC is
still ineffective to search for di-tau events from the direct stau productions
[32]. These limits are sufficiently weak. Fig. 5 does not change even if they
are imposed.
| $m_{\tilde{\ell}}$ | $M_{1}$ | $\mu$ | $m_{\tilde{e}_{1}}$ | $m_{\tilde{e}_{2}}$ | $m_{\tilde{\mu}_{1}}$ | $m_{\tilde{\mu}_{2}}$ | $m_{\tilde{\tau}_{1}}$ | $m_{\tilde{\tau}_{2}}$ | $m_{\tilde{\chi}^{0}_{1}}$ | $\Delta a_{\mu}$
---|---|---|---|---|---|---|---|---|---|---|---
A | 300 | 200 | 756 | 303 | 304 | 298 | 309 | 199 | 380 | 199 | 16.1
$\rm{A^{\prime}}$ | 300 | 200 | 699 | 303 | 304 | 299 | 308 | 209 | 375 | 199 | 14.6
B | 470 | 250 | 1680 | 472 | 472 | 465 | 479 | 329 | 581 | 250 | 10.2
C | 340 | 160 | 1138 | 343 | 343 | 336 | 350 | 199 | 442 | 160 | 18.0
Table 2: Model parameters and mass spectra at several model points in Fig. 5.
The masses are in units of GeV, and the muon $g-2$ is scaled by $10^{-10}$.
Here, $m_{\tilde{\ell}}$ denotes $m_{\tilde{\ell}_{L}}=m_{\tilde{\ell}_{R}}$,
and $\tan\beta=40$ is set.
In Tab. 2, superparticle mass spectra are listed for several points in Fig. 5.
The lightest neutralino mass is close to the Bino mass, since $\mu$ is very
large. Given $m_{\tilde{\ell}}\equiv
m_{\tilde{\ell}_{L}}=m_{\tilde{\ell}_{R}}$, the slepton masses are
hierarchical except for those of the selectrons due to a large left-right
mixing. The lightest stau mass is closest to the neutralino mass among the
sleptons. In Fig. 5, the point A (A’) is around the upper side of the contour
of the muon $g-2$. In the vicinity of this side, $m_{\tilde{\tau}_{1}}$ has
the closest value to $m_{\tilde{\chi}^{0}_{1}}$ in the allowed range. If
$m_{\tilde{\tau}_{1}}>m_{\tilde{\chi}^{0}_{1}}$ is imposed, it satisfies
$m_{\tilde{\tau}_{1}}=m_{\tilde{\chi}^{0}_{1}}$ to maximize the SUSY
contributions to the muon $g-2$ around the upper side of the contour. Above
it, the regions are excluded by the long-lived stau search, or the muon $g-2$
becomes too small. On the other hand, if the mass difference between the stau
and the neutralino is assumed to be larger than $\delta M$, they satisfy
$m_{\tilde{\tau}_{1}}=m_{\tilde{\chi}^{0}_{1}}+\delta M$ in the vicinity of
the upper side. At the point A’, $\delta M=10\,{\rm GeV}$ is imposed, and
$m_{\tilde{\tau}_{1}}$ is larger than $m_{\tilde{\chi}^{0}_{1}}$ by $10\,{\rm
GeV}$. In contrast, the points B and C are away from this region. B (C) is
close to the maximal end point of the lightest smuon mass which explains the
muon $g-2$ at the $2\sigma$ ($1\sigma$) level. Here, the left-right mixing is
determined by the vacuum stability condition of the stau.
Figure 7: The production cross section of the lightest smuon at ILC with
$\sqrt{s}=1\,{\rm TeV}$. The beams are assumed to be polarized at 80%
$(e^{-})$ and 30% $(e^{+})$. The left-right mixing is maximized under the
vacuum stability condition with $\tan\beta=40$. The $1\sigma$ ($2\sigma$)
regions of the muon $g-2$ are also shown by the orange (yellow) bands. The
parameters are $M_{1}=200\,{\rm GeV}$ and $M_{\rm soft}=10\,{\rm TeV}$.
Linear colliders of $e^{+}e^{-}$ (ILC) are very useful (see e.g., Ref. [33,
34]). They can provide rich informations of the models. Moreover, they are
superior to LHC when the slepton masses are close to that of the lightest
neutralino [35]. In the orange/yellow regions of Fig. 5, selectrons, smuons
and staus can be produced at the linear colliders. In Fig. 7, production cross
sections of the lightest smuon are shown for the left- and right-handed smuon
masses. Here, it is assumed that the collision energy is $\sqrt{s}=1\,{\rm
TeV}$ at ILC, and the beams are polarized at 80% for the electron and 30% for
the positron. The left-right mixing is maximized under the vacuum stability
condition. The other parameters are $M_{1}=200\,{\rm GeV}$ and $\tan\beta=40$.
It is found that the cross section is larger than 20 (1)$\,{\rm fb^{-1}}$ for
the $1\sigma$ ($2\sigma$) parameter region of the muon $g-2$. The estimation
is based on Ref. [36].888 The cross section is calculated at the leading order
in Fig. 7. It can be enhanced by $\sim 10\%$ near the mass threshold [36]. It
is expected that the smuons can be discovered in almost all the parameter
region that are kinematically allowed (cf. [34]).
It is possible to measure masses of the smuon and the neutralino from event
distributions and the cross section [37, 13]. Also, the chirality structure of
the smuons could be determined by the beam polarization. We assume
$(P_{e^{-}},P_{e^{+}})=(+0.8,-0.3)$ in the left panel of Fig. 7, and
$(-0.8,+0.3)$ in the right panel. The cross section of a chiral smuon is
sensitive to the polarization, because the productions proceed by the
s-channel $\gamma/Z$ exchanges. On the other hand, when
$m_{\tilde{\mu}_{L}}=m_{\tilde{\mu}_{R}}$, the cross section is insensitive to
the beam polarization. Since the left-right mixing is large, the left- and
right-handed smuons are maximally mixed with each other. Both of them can be
produced, and the mass difference is expected to be measured, for instance, by
measuring threshold productions. At the model points in Tab. 2, the difference
is $\gtrsim 10\,{\rm GeV}$, which is much larger than the uncertainty of the
mass measurement at ILC, $\mathop{\mathcal{O}}(10-100)\,{\rm MeV}$ [33]. It is
emphasized that the smuon productions are clean and direct signatures of the
SUSY contributions to the muon $g-2$.
Figure 8: The production cross section of $\tilde{e}^{+}_{R}\tilde{e}^{-}_{L}$
and $\tilde{e}^{+}_{L}\tilde{e}^{-}_{R}$ at ILC with $\sqrt{s}=1\,{\rm TeV}$.
The parameters are same as Fig. 7, but $M_{1}$ is varied with
$m_{\tilde{e}_{L}}=m_{\tilde{e}_{R}}\equiv m_{\tilde{e}}$. The corrections to
the Bino coupling due to decoupling of heavy superparticles are included
(discarded) in the black solid (dashed) lines.
The selectron productions proceed not only by the s-channel $\gamma/Z$
exchanges, but also by the t-channel Bino exchange. The latter contribution
can enhance the cross section and provide individual information in addition
to mass measurements of the selectron and the neutralino. The productions,
$e^{+}_{L}e^{-}_{L}\to\tilde{e}^{+}_{R}\tilde{e}^{-}_{L}$ and
$e^{+}_{R}e^{-}_{R}\to\tilde{e}^{+}_{L}\tilde{e}^{-}_{R}$, proceed by the
t-channel. In Fig. 8, their total cross sections are shown by the black solid
lines. Here, the parameters are $m_{\tilde{e}_{L}}=m_{\tilde{e}_{R}}$, and the
beam polarization is 80% for the electron and 30% for the positron. It is
found that the cross sections are $\mathop{\mathcal{O}}(10)\,{\rm fb}$ in the
muon $g-2$ parameter regions. They are larger than those of the smuon in Fig.
7. Importantly, this channel is useful to measure the Bino–electron–selectron
couplings [13, 14]. They are deviated from the U(1)Y gauge coupling constant
by decoupling heavy superparticles, as discussed in Sec. 3.1. In Fig. 8, the
cross sections are estimated with (without) $\delta\tilde{g}_{L}$ and
$\delta\tilde{g}_{R}$ (see Eqs. (11) and (12)). The results are shown by the
solid (dashed) lines. It is found that the corrections enhance the cross
section by $8-10\%$ for $M_{\rm susy}=10\,{\rm TeV}$. This is measurable at
ILC [13, 14].
Cross sections of $e^{+}_{L}e^{-}_{R}\to\tilde{e}^{+}_{R}\tilde{e}^{-}_{R}$
and $e^{+}_{R}e^{-}_{L}\to\tilde{e}^{+}_{L}\tilde{e}^{-}_{L}$ also include the
t-channel contribution of the Bino (see e.g., Ref. [36]). In particular, the
former cross section is enhanced well compared to those solely by the
s-channel $\gamma/Z$ exchanges. It can be $\gtrsim 100\,{\rm fb}$ when $M_{1}$
is relatively small in the muon $g-2$ parameter regions. However, in the
latter process, the t-channel contribution interferes destructively with the
s-channel contribution. On the other hand, the cross section of
$e^{+}_{L}e^{-}_{R}\to\tilde{e}^{+}_{R}\tilde{e}^{-}_{R}$ differs by $5-6\%$
between the cases with and without $\delta\tilde{g}_{L}$ and
$\delta\tilde{g}_{R}$. This is smaller than the above channels.
The stau productions proceed similarly to the smuon at ILC. The stau mass
parameters can be determined [38, 13, 39]. The produced staus decay into the
lightest (Bino-like) neutralino and a tau lepton. The tau polarization can be
measured from the energy spectra of tau hadronic decays [38]. If productions
of both the heavy and light staus are kinematically allowed, all the
components of the stau mass matrix can be measured. In particular, the left-
right mixing is determined.999 Strictly, it is difficult to determine $\mu$
and $A_{\tau}$ separately by the stau productions, because the lightest
neutralino is almost composed of the Bino [13, 39].
Let us discuss the dark matter in the present model. The lightest neutralino
is a candidate of the dark matter. The neutralino relic abundance becomes
consistent with the measured cold dark matter abundance [40, 41] by the stau
co-annihilation. The mass difference between the stau and the neutralino is
required to be $5-10\,{\rm GeV}$. This corresponds to the model points in the
vicinity of the upper side of the muon $g-2$ contours in Fig. 5, including the
points A or A’. This region has not been excluded by the studies of the di-
lepton signature at LHC [29, 30]. Masses of the selectrons and the smuons are
too close to that of the neutralino for LHC. ILC is superior to study the
region [33, 34]. Selectrons and smuons as well as staus can be produced, as
discussed above.
Finally, let us mention the branching ratios of the Higgs boson decays. When
the staus are light, and their left-right mixing is large, the branching ratio
of the di-photon channel can be enhanced sizably [42]. We have estimated it in
the muon $g-2$ parameter regions of Fig. 5. In the $1\sigma$ parameter region,
the ratio increases by $10-40\%$ compared to the SM prediction, while it is
enhanced by $5-10\%$ in the $2\sigma$ region. The latter is comparable to the
sensitivity of the high-luminosity LHC with $\int\mathcal{L}=3000\,{\rm
fb^{-1}}$ [43, 44]. On the other hand, the ratio of the Higgs boson decaying
to the di-muon is almost unchanged. Even when the radiative corrections to the
Higgs boson coupling with the muon is very large, i.e., $\Delta_{\mu}\gg 1$,
the ratio does not change, because the tree level (Yukawa) coupling decreases
by $\Delta_{\mu}$, and the sum of the tree and radiative contributions is not
varied.
### 4.2 Non-universal slepton mass
Figure 9: The lower bound on the stau mass to satisfy the vacuum stability
bound of the stau–Higgs potential. In the orange (yellow) region, the SUSY
contribution to the muon $g-2$ can explain the discrepancy (1) at the
$1\sigma$ $(2\sigma)$ level. The parameters are
$m_{\tilde{\mu}_{L}}=m_{\tilde{\mu}_{R}}=M_{1}$,
$m_{\tilde{\tau}_{L}}=m_{\tilde{\tau}_{R}}$, $\tan\beta=40$ and $M_{\rm
soft}=30\,{\rm TeV}$.
In this section, we discuss the case without slepton mass universality. The
slepton soft SUSY-breaking masses satisfy
$\displaystyle
m_{\tilde{e}_{L}}=m_{\tilde{\mu}_{L}}<m_{\tilde{\tau}_{L}},~{}~{}~{}m_{\tilde{e}_{R}}=m_{\tilde{\mu}_{R}}<m_{\tilde{\tau}_{R}}.$
(17)
As discussed in Sec. 3.2, the vacuum stability bound of the stau–Higgs
potential is relaxed by heavy staus, and the smuon masses are allowed to be
larger in order to solve the muon $g-2$ anomaly. In other words, the stau
masses are required to be large relative to the smuons, when the smuons are
heavy. In Fig. 9, the lower bound on the stau mass is shown. In the orange
(yellow) regions, the muon $g-2$ anomaly is solved at the $1\sigma$
($2\sigma$) level, while the vacuum stability constraint is avoided. The
ratios of the stau and smuon masses are defined as
$\displaystyle
R_{L}\equiv\frac{m_{\tilde{\tau}_{L}}}{m_{\tilde{\mu}_{L}}},~{}~{}~{}R_{R}\equiv\frac{m_{\tilde{\tau}_{R}}}{m_{\tilde{\mu}_{R}}}.$
(18)
The parameters are $M_{1}=m_{\tilde{\mu}_{L}}=m_{\tilde{\mu}_{R}}$,
$m_{\tilde{\tau}_{L}}=m_{\tilde{\tau}_{R}}$, $\tan\beta=40$ and $M_{\rm
soft}=30\,{\rm TeV}$. The vertical axis is $R\equiv R_{L}=R_{R}$. If the staus
are decoupled, the left-right mixing is bounded either by Eq. (15) or the
vacuum stability bound of the smuon–Higgs potential, as discussed in Sec. 3.2.
The smuon masses can be $1.4\,{\rm TeV}$ $(1.9\,{\rm TeV})$ for the $1\sigma$
$(2\sigma)$ level of the muon $g-2$. Even though it becomes difficult to
produce such sleptons at LHC or ILC, non-universal slepton masses generically
cause problems of too large lepton flavor violations (LFV) and CP violations
(CPV).
LFV and CPV are sensitive to off-diagonal (generation mixing) components of
the slepton mass matrices. They are suppressed if the slepton mass matrices
are universal among the flavors, such as in the gauge mediated SUSY-breaking
models. However, even when the mass matrices are diagonal in the model basis,
sizable FCNC and CPV are generically induced as long as the diagonal
components are not equal to each others. In fact, a lot of models have been
proposed to explain the SM Yukawa couplings, and many of them predict non-
diagonal Yukawa matrices in the model basis. Even if the sfermion mass
matrices are diagonal in this basis, off-diagonal components are generated in
the fermion mass eigenstate basis, which are obtained by rotating the mass
matrices of the model basis with unitary matrices. In other words, the super
GIM mechanism does not work generically unless the sfermion mass matrices are
universal.
When the muon $g-2$ anomaly is solved by the SUSY contributions, LFV and
lepton electric dipole moments (EDM) are generically sizable. In this section,
the following setup is considered. In the model basis, the slepton mass
matrices are diagonal among the flavors,
$\displaystyle(m_{\tilde{\ell}_{L}}^{2})_{ij}$
$\displaystyle=\text{diag}(m_{\tilde{e}_{L}}^{2},m_{\tilde{\mu}_{L}}^{2},m_{\tilde{\tau}_{L}}^{2}),$
(19) $\displaystyle(m_{\tilde{\ell}_{R}}^{2})_{ij}$
$\displaystyle=\text{diag}(m_{\tilde{e}_{R}}^{2},m_{\tilde{\mu}_{R}}^{2},m_{\tilde{\tau}_{R}}^{2}).$
(20)
Moreover, $m_{\tilde{e}_{L}}=m_{\tilde{\mu}_{L}}$ and
$m_{\tilde{e}_{R}}=m_{\tilde{\mu}_{R}}$ are imposed. Otherwise, LFV becomes
too large, as shown later. The Yukawa matrices are generally non-diagonal in
the model basis. The mass eigenstate basis of the charged leptons is obtained
by the left- and right-handed unitary matrices, $U_{L},U_{R}$, as
$\displaystyle U_{R}\,M_{\ell}\,U_{L}^{\dagger}\equiv
U_{R}\,Y_{\ell}\,v_{d}\,U_{L}^{\dagger}={\rm
diag}\left(\frac{m_{e}}{1+\Delta_{e}},~{}\frac{m_{\mu}}{1+\Delta_{\mu}},~{}\frac{m_{\tau}}{1+\Delta_{\tau}}\right),$
(21)
where $v_{d}$ is the down-type Higgs VEV, and $\Delta_{\ell}$ is given by Eq.
(7). The unitary matrices are generally represented as
$U_{L,R}=\exp\left[\begin{pmatrix}0&(\delta_{L,R})_{12}&(\delta_{L,R})_{13}\\\
-(\delta_{L,R})^{*}_{12}&0&(\delta_{L,R})_{23}\\\
-(\delta_{L,R})^{*}_{13}&-(\delta_{L,R})^{*}_{23}&0\\\ \end{pmatrix}\right].$
(22)
Non-vanishing mixings, $(\delta_{L,R})_{ij}$, induce LFV and EDM, as long as
the slepton masses are non-universal.101010 Similar setup has been studied in
Ref. [45], where squarks of the first two generations are light. Note that
quark FCNC and CPV are suppressed by heavy colored superparticles in this
paper.
The magnetic dipole contributions to LFV are represented by the following
effective Lagrangian,
$\displaystyle\mathcal{L}_{\rm
eff}=e\frac{m_{\ell_{j}}}{2}\bar{\ell}_{i}\sigma_{\mu\nu}\left(A^{L}_{ij}P_{L}+A^{R}_{ij}P_{R}\right)\ell_{j}F^{\mu\nu}+{\rm
h.c.},$ (23)
where $i,j$ are flavor indices. Contributions to the other higher dimensional
operators are subdominant in this paper. The Wilson coefficients, $A^{L}_{ij}$
and $A^{R}_{ij}$, are dominated by the Bino–slepton contributions, similarly
to the muon $g-2$. In the mass insertion approximation, they are estimated as
(cf. [46])
$\displaystyle A^{L}_{ij}$ $\displaystyle=(1+\delta^{\rm
2loop})\frac{\alpha_{Y}}{8\pi}\frac{M_{1}\mu\tan\beta}{m_{\ell_{j}}}\sum_{a,b=1,2,3}\Big{[}U_{R}\Big{]}_{ib}\Big{[}M_{\ell}\Big{]}_{ba}\Big{[}U_{L}^{\dagger}\Big{]}_{aj}\,F_{a,b},$
(24) $\displaystyle A^{R}_{ij}$ $\displaystyle=(1+\delta^{\rm
2loop})\frac{\alpha_{Y}}{8\pi}\frac{M_{1}\mu\tan\beta}{m_{\ell_{j}}}\sum_{a,b=1,2,3}\Big{[}U_{L}\Big{]}_{ia}\Big{[}M_{\ell}^{\dagger}\Big{]}_{ab}\Big{[}U_{R}^{\dagger}\Big{]}_{bj}\,F_{a,b}.$
(25)
The two loop factor $\delta^{\rm 2loop}$ is found in Eq. (9). The loop
function $F_{a,b}$ is defined as
$\displaystyle
F_{a,b}=\frac{1}{m_{\tilde{\ell}_{La}}^{2}m_{\tilde{\ell}_{Rb}}^{2}}f_{N}\left(\frac{m_{\tilde{\ell}_{La}}^{2}}{|M_{1}|^{2}},\frac{m_{\tilde{\ell}_{Rb}}^{2}}{|M_{1}|^{2}}\right).$
(26)
The muon FCNC’s are the most sensitive to the non-universal slepton mass
matrices and the non-diagonal Yukawa matrix. The decay rate of $\mu\to
e\gamma$ is represented as
$\displaystyle\Gamma(\mu\to
e\gamma)=\frac{\alpha}{4}m_{\mu}^{5}\left(\left|A^{L}_{12}\right|^{2}+\left|A^{R}_{12}\right|^{2}\right),$
(27)
where the unitarity matrices in $A^{R}_{12}$ are expanded as
$\displaystyle\sum_{a,b=1,2,3}\Big{[}U_{L}\Big{]}_{1a}\Big{[}M_{\ell}^{\dagger}\Big{]}_{ab}\Big{[}U_{R}^{\dagger}\Big{]}_{b2}\,F_{a,b}$
$\displaystyle=-\frac{m_{\mu}}{1+\Delta_{\mu}}(\delta_{L})_{12}\left(F_{1,2}-F_{2,2}\right)$
(28)
$\displaystyle~{}~{}~{}+\frac{m_{\tau}}{1+\Delta_{\tau}}(\delta_{L})_{13}(\delta_{R})_{23}^{*}\left(F_{1,2}-F_{1,3}-F_{3,2}+F_{3,3}\right),$
at the leading order of $(\delta_{L})_{ij}$, $(\delta_{R})_{ij}$ and
$m_{\mu}/m_{\tau}$. Here and hereafter, $m_{e}=0$ is set, for simplicity.
Similarly, $A^{L}_{12}$ is obtained by replacing $L\leftrightarrow R$. In the
last term, $F_{1,2}$ is dominant when the staus are heavy, i.e., $F_{1,2}\gg
F_{1,3},F_{3,2},F_{3,3}$ for $m_{\tilde{e}_{L}},m_{\tilde{\mu}_{L}}\ll
m_{\tilde{\tau}_{L}}$ and $m_{\tilde{e}_{R}},m_{\tilde{\mu}_{R}}\ll
m_{\tilde{\tau}_{R}}$. On the other hand, the right-hand side vanishes when
the slepton masses are universal, as expected from the super GIM mechanism. In
particular, when the sleptons are degenerate in the first two generations, the
first term becomes zero due to $F_{1,2}=F_{2,2}$. Otherwise, the muon LFV is
induced at the order of $(\delta_{L})_{12}$.
The above rate is compared with the SUSY contribution to the muon $g-2$. In
the non-universal slepton mass spectrum, it is represented by $A^{L}_{ij}$ and
$A^{R}_{ij}$ as
$\displaystyle a_{\mu}({\rm SUSY})$ $\displaystyle=m_{\mu}^{2}\,{\rm
Re}\left[A^{L}_{22}+A^{R}_{22}\right]$ $\displaystyle=(1+\delta^{\rm
2loop})\frac{\alpha_{Y}}{4\pi}m_{\mu}M_{1}\mu\tan\beta\left[\frac{m_{\mu}}{1+\Delta_{\mu}}F_{2,2}+\kappa\right].$
(29)
This is same as Eq. (4) up to a correction $\kappa$, which is represented as
$\displaystyle\kappa=\frac{m_{\tau}}{1+\Delta_{\tau}}{\rm
Re}\left[(\delta_{L})_{23}(\delta_{R})_{23}^{*}\right]\left(F_{2,2}-F_{2,3}-F_{3,2}+F_{3,3}\right)+\cdots.$
(30)
Here, the omitted terms are suppressed by $(\delta_{L})_{ij}$,
$(\delta_{R})_{ij}$ or $m_{\mu}/m_{\tau}$. If the slepton mass matrices are
universal, $\kappa$ vanishes.
It is noticed that Eqs. (27) and (29) are tightly correlated to each other.
For the mass spectrum (17) with $R_{L},R_{R}\gg 1$, the ratio is
$\displaystyle\frac{{\rm Br}(\mu\to e\gamma)}{a_{\mu}({\rm
SUSY})^{2}}\simeq\frac{1}{\Gamma_{\rm tot}}\frac{\alpha
m_{\mu}}{16}\,\big{|}\delta_{13}\delta_{23}\big{|}^{2}\left(\frac{m_{\tau}}{m_{\mu}}\frac{1+\Delta_{\mu}}{1+\Delta_{\tau}}\right)^{2},$
(31)
where $\Gamma_{\rm tot}$ is the total decay rate of muon, and the mixing is
defined as
$\displaystyle\big{|}\delta_{13}\delta_{23}\big{|}^{2}\equiv\big{|}(\delta_{R})_{13}(\delta_{L})_{23}\big{|}^{2}+\big{|}(\delta_{L})_{13}(\delta_{R})_{23}\big{|}^{2}.$
(32)
It is independent of superparticle mass spectra except through $\Delta_{\mu}$
and $\Delta_{\tau}$. Thus, when the muon $g-2$ discrepancy (1) is explained by
the SUSY contributions, $\mu\to e\gamma$ is induced by the non-universal
slepton mass sizably. It is important that the decay is not suppressed by
heavy slepton masses, for instance, $m_{\tilde{\mu}_{1}}=1.4\,{\rm TeV}$ or
$1.9\,{\rm TeV}$ in Tab. 1, for given SUSY contributions to the muon $g-2$.
Figure 10: Contours of ${\rm Br}(\mu\to e\gamma)$ (left) and $d_{e}$ (right)
for $a_{\mu}({\rm SUSY})=10^{-9}$ with $M_{\rm soft}=30\,{\rm TeV}$. In the
left panel, the upper black solid line is the current bound by MEG, ${\rm
Br}(\mu\to e\gamma)<5.7\cdot 10^{-13}$ at 90% CL [47]. The lower one is the
sensitivity of the MEG upgrade, ${\rm Br}(\mu\to e\gamma)=6\cdot 10^{-14}$
[48]. The blue dashed line is the Mu3e sensitivity, ${\rm Br}(\mu\to
e\bar{e}e)=10^{-16}$ [49]. The $\mu-e$ conversion is expected to probe down to
$R_{\mu e}=3\cdot 10^{-17}$ (upper red dotted) by COMET/Mu2e [50, 51] and
$2\cdot 10^{-19}$ (lower red dotted) by PRISM/PRIME [52]. In the right panel,
the black solid line is the current bound by using YbF, $|d_{e}|<1.05\cdot
10^{-27}\,e{\rm cm}$ at 90% CL [53]. The sensitivity is planned to be
improved: $|d_{e}|=10^{-29}\,e{\rm cm}$ by Fr or ThO [54, 55, 56], and
$10^{-30}\,e{\rm cm}$ by YbF or WN [57, 58]. Here,
$m_{\tilde{\ell}_{L}}=m_{\tilde{\ell}_{R}}$ is assumed.
In the left panel of Fig. 10, contours of ${\rm Br}(\mu\to e\gamma)$ are
shown. Here, $m_{\tilde{\ell}_{L}}=m_{\tilde{\ell}_{R}}$ is assumed. The SUSY
contributions to the muon $g-2$ is fixed to be $a_{\mu}({\rm SUSY})=1\times
10^{-9}$ with $M_{\rm soft}=30\,{\rm TeV}$. The small corrections are taken at
$M_{1}=m_{\tilde{\mu}}=400\,{\rm GeV}$ and $\tan\beta=40$ as a reference,
though the result is almost independent of them. In the figure, the contours
correspond to the current limit and future sensitivities of experiments,
* •
the current limit of the MEG experiment, ${\rm Br}(\mu\to e\gamma)<5.7\times
10^{-13}$ at 90% CL [47] (upper black solid line in the figure).
* •
the sensitivity of the MEG upgrade, ${\rm Br}(\mu\to e\gamma)=6\times
10^{-14}$ [48] (lower black solid line).
* •
the sensitivity of the Mu3e experiment, ${\rm Br}(\mu\to e\bar{e}e)=10^{-16}$
at Phase II [49] (blue dashed line).
* •
the sensitivity of the COMET experiment, $R_{\mu e}=3\times 10^{-17}$ at Phase
II [50] (upper red dotted line). The Mu2e experiment has a similar sensitivity
[51].
* •
the proposal of the PRISM/PRIME project, $R_{\mu e}=2\times 10^{-19}$ [52]
(lower red dotted line).
Note that $\mu\to e\bar{e}e$ and $\mu-e$ conversion experiments have better
sensitivity in future than those of $\mu\to e\gamma$. In particular, the
latter experiment has low (accidental) backgrounds. On the other hand, the
current constraint and future sensitivities of the tau LFV’s are weaker than
those of the muon, though they are also induced by the non-universal slepton
masses with finite $(\delta_{L})_{23}$ and $(\delta_{R})_{23}$.
From the figure, it is found that the LFV decay rate increases rapidly as the
staus become heavier than the smuons, $R>1$. When the staus are decoupled,
$R\gg 1$, the mixing $|\delta_{13}\delta_{23}|$ in Eq. (32) is limited to be
smaller than $3\times 10^{-6}$ by MEG for $a_{\mu}({\rm SUSY})=1\times
10^{-9}$. For instance, when the smuon and selectron masses are larger than
$1\,{\rm TeV}$, the stau are required to be heavier than $7\,{\rm TeV}$
($R>7$) to avoid the vacuum stability constraint, according to Fig. 9. If the
lepton Yukawa matrix is related to the quark sector, e.g., by the GUT
relation, it is naively expected to be $|\delta_{13}\delta_{23}|\sim
V_{ub}V_{cb}\sim 10^{-4}$. This already exceeds the above limit. Thus, non-
universal slepton mass spectra are tightly constrained by LFV. In future, if
sleptons are neither discovered at LHC nor ILC, the model is expected to be
probed by LFV. Otherwise, the SM Yukawa matrices are tightly limited in the
model basis, or when flavor models are constructed.
The flavor off-diagonal components of the Yukawa matrices are sources of the
CP violations, when the slepton mass matrices are non-universal. Similarly to
the correction $\kappa$ of the muon $g-2$ in Eq. (29), the electron EDM is
induced as
$\displaystyle\frac{d_{e}}{e}$ $\displaystyle=\frac{m_{e}}{2}\,{\rm
Im}\left[A^{L}_{11}-A^{R}_{11}\right]$ $\displaystyle=(1+\delta^{\rm
2loop})\frac{\alpha_{Y}}{8\pi}M_{1}\mu\tan\beta$
$\displaystyle~{}~{}~{}\times\bigg{[}\frac{m_{\mu}}{1+\Delta_{\mu}}{\rm
Im}\left[(\delta_{R})_{12}(\delta_{L})_{12}^{*}\right]\left(F_{1,1}-F_{1,2}-F_{2,1}+F_{2,2}\right)$
$\displaystyle~{}~{}~{}~{}~{}~{}+\frac{m_{\tau}}{1+\Delta_{\tau}}{\rm
Im}\left[(\delta_{R})_{13}(\delta_{L})_{13}^{*}\right]\left(F_{1,1}-F_{1,3}-F_{3,1}+F_{3,3}\right)+\cdots\bigg{]}.$
(33)
Here, we assume $\arg(M_{1}\mu\tan\beta)=0$. The omitted terms are suppressed
by orders of $(\delta_{L})_{ij}$, $(\delta_{R})_{ij}$ or $m_{\mu}/m_{\tau}$.
In the last term, $F_{1,1}\gg F_{1,3},F_{3,1},F_{3,3}$ is obtained when the
staus are heavy. On the other hand, the right hand side vanishes when the
slepton masses are universal, because the complex phases can be rotated away.
Comparing Eq. (33) with Eq. (29), we obtain
$\displaystyle\frac{d_{e}/e}{a_{\mu}({\rm
SUSY})}\simeq\frac{1}{2m_{\mu}}{\rm{Im}}[(\delta_{R})_{13}(\delta_{L})^{\ast}_{13}]\,\frac{m_{\tau}}{m_{\mu}}\frac{1+\Delta_{\mu}}{1+\Delta_{\tau}},$
(34)
for the mass spectrum (17) with $R_{L},R_{R}\gg 1$. It is independent of
superparticle mass spectra except through $\Delta_{\mu}$ and $\Delta_{\tau}$.
Thus, if the muon $g-2$ anomaly is solved by the SUSY contributions, EDM
becomes sizable by the non-universal slepton mass. Similarly to LFV, this is
valid even for large selectron and smuon masses.
In the right panel of Fig. 10, contours of the electron EDM are shown. Here,
the SUSY contributions to the muon $g-2$ is fixed to be $a_{\mu}({\rm
SUSY})=1\times 10^{-9}$ with $M_{\rm soft}=30\,{\rm TeV}$. The result is
almost independent of superparticle mass spectra except for small corrections,
$\Delta_{\mu}$ and $\Delta_{\tau}$. The following data are used,
* •
the current limit with the YbF molecule, $|d_{e}|<1.05\times 10^{-27}\,e{\rm
cm}$ at 90% CL [53] (solid line in the figure).
* •
the sensitivity with the Fr atom, $|d_{e}|=10^{-29}\,e{\rm cm}$ [54] (dashed
line). The experiment with the ThO molecule could have a similar sensitivity
by accumulating data, $|d_{e}|=1\times 10^{-28}/\sqrt{({\rm day})}\,e{\rm cm}$
[55, 56].
* •
the sensitivity with the YbF molecule, $|d_{e}|=10^{-30}\,e{\rm cm}$ [57]
(dotted line). The experiment with the WN ion can probe down to,
$|d_{e}|=10^{-30}\,e{\rm cm}/{\rm day}$, where the systematic limit is at the
level of $10^{-31}\,e{\rm cm}$ [58].
From the figure, it is found that EDM is sensitive to
${\rm{Im}}[(\delta_{R})_{13}(\delta_{L})^{*}_{13}]$ when the staus are heavier
than the selectrons. The current experimental limit puts a constraint,
${\rm{Im}}[(\delta_{R})_{13}(\delta_{L})^{*}_{13}]<6\times 10^{-7}$ for $R\gg
1$ and $a_{\mu}({\rm SUSY})=1\times 10^{-9}$. The sensitivity will be improved
very well. The mixing will be able to be probed at the level of $10^{-10}$.
Thus, if sleptons are neither discovered at LHC nor ILC in future, the model
can be sensitively probed by EDM as well as LFV. If no signal will be
observed, the CP violating phase must be suppressed very tightly in order to
explain the muon $g-2$ anomaly by the SUSY contributions.
## 5 Conclusion
SUSY is one of the most motivated candidates of the new physics. If the muon
$g-2$ anomaly is solved by SUSY, some of the superparticles are relatively
light. In this paper, we focused on the Bino–smuon contribution to the muon
$g-2$. It is enhanced not only by $\tan\beta$ but also by large $\mu$.
Consequently, it was shown that the smuon masses can be as large as $\sim
1\,{\rm TeV}$. We examined the phenomenology of the models, in which only the
superparticles that are relevant for the muon $g-2$ are light.
The analyses were categorized by the slepton mass spectrum. When the mass
spectrum is universal among the flavors, it was found that the vacuum
stability of the stau–Higgs potential restricts the smuon masses tightly. They
are predicted to be within $330~{}(460)\,{\rm GeV}$ at the
$1\sigma~{}(2\sigma)$ level of the muon $g-2$. It was shown that part of the
parameter region is already excluded by LHC, and argued that such sleptons are
expected to be studied at LHC or ILC in future.
If the staus are (much) heavier than the smuons, the vacuum stability bound of
the staus is relaxed. In this mass spectrum, the smuon masses are limited by
the vacuum stability condition of the smuon–Higgs potential. It was found that
they are less than $1.4~{}(1.9)\,{\rm TeV}$ at the $1\sigma~{}(2\sigma)$ level
of the muon $g-2$. Such slepton masses exceed the LHC/ILC reach. Instead, the
non-universal slepton mass spectrum generically predicts too large LFV and
EDM. They originate in the non-diagonal SM Yukawa matrices in the model basis.
Although the prediction depends on the flavor models, LFV and/or EDM is very
likely to be sizable in wide models. If no sleptons are discovered at LHC nor
ILC, the model is expected to be probed by LFV and EDM sensitively.
So far, $\tan\beta=40$ was chosen in the numerical analyses. As mentioned in
Sec. 2, the observables depend on $m_{\tilde{\ell}_{LR}}^{2}$, and $\tan\beta$
contributes to them only through it. Thus, the results are independent of a
choice of $\tan\beta$ up to a correction to the vacuum stability condition
(see Sec. 3.2). Since the correction is small, the above conclusion does not
change, even if $\tan\beta$ is varied, for instance, from 20 to 70.
In conclusion, we studied the SUSY models that explain the muon $g-2$
discrepancy with the Bino–smuon contributions. It was shown that they are
expected to be probed by LHC/ILC and LFV/EDM complementarily in future.
## Acknowledgements
This work was supported by JSPS KAKENHI Grant No. 23740172 (M.E.), 22244021
(K.H.) and 25-10486 (T.Y.). The work of T.K. is partially supported by Global
COE Program “the Physical Sciences Frontier”, MEXT, Japan. The work of T.Y. is
supported in part by a JSPS Research Fellowship for Young Scientists, and
supported by an Advanced Leading Graduate Course for Photon Science grant.
## Note Added
While we are finalizing the paper, the article [59] was submitted to arXiv.
The article studied non-decoupling two-loop contributions from heavy
sfermions, which partially overlaps with the discussion on the correction to
the muon $g-2$ via the Bino–muon–smuon coupling in Sec. 3.1.
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|
arxiv-papers
| 2013-09-12T09:07:44 |
2024-09-04T02:49:50.804074
|
{
"license": "Public Domain",
"authors": "Motoi Endo, Koichi Hamaguchi, Teppei Kitahara, and Takahiro Yoshinaga",
"submitter": "Takahiro Yoshinaga",
"url": "https://arxiv.org/abs/1309.3065"
}
|
1309.3095
|
# Propagation of an arbitrary vortex pair through an astigmatic optical system
and determination of its net topological charge
Salla Gangi Reddy Corresponding author: [email protected] Shashi Prabhakar
[email protected] Aadhi A [email protected] J. Banerji [email protected] R. P.
Singh [email protected] Physical Research Laboratory, Navrangpura,
Ahmedabad, India - 380009
###### Abstract
We embed a pair of vortices with different topological charges in a Gaussian
beam and study its evolution through an astigmatic optical system, a tilted
lens. The propagation dynamics is explained by a closed-form analytical
expression. Furthermore, we show that a careful examination of the intensity
distribution at a predicted position past the lens can provide us with the net
charge present in the beam. To the best of our knowledge, our method is the
first non-interferometric technique to measure the net charge of an arbitrary
vortex pair. Our theoretical results are well supported by experimental
observations.
###### pacs:
050:0050, 050.4865.
## I Introduction
Optical vortices have drawn considerable attention in science and engineering
due to their dark core and helical wave front. An optical vortex of order $l$
centered at the origin ($r=0$) has a field distribution of the form
$E(r)\exp(il\phi)$. The distribution is such that the field intensity tends to
zero as $r\rightarrow 0$ whereas the phase shift in one cycle around the
origin is $2l\pi$ where $l$ is an integer. The azimuthal mode index $l$, also
called topological charge of the vortex, has a physical meaning in that the
vortex carries an orbital angular momentum (OAM) of $l\hbar$ per photon allen
. This angular momentum can be imparted to microscopic particles in order to
manipulate them optically particle ; particle1 ; grier . In recent years, the
OAM of light has also found applications in classical siddu as well as
quantum communication OAM ; torner . These applications have led to
considerable interest in the generation and study of optical vortices both in
free space free ; spiral ; cgh and in guided media guide ; guide1 .
The propagation of a pair of vortices has gained a lot of interest in research
since the last decade dipole1 ; dipole2 ; dipole3 ; dipole4 ; dipole5 ;
dipole6 . Indebetouw studied the propagation of an array of vortices through
free space and showed that the relative separation between the vortices is
invariant during the propagation in the case of same type (sign) of charges
whereas they will attract and annihilate each other in the case of oppositely
charged vortices dipole1 . Chen and Roux have studied the annihilation of
dipole vortices during their propagation. They found that the background phase
function at a point where two dipoles annihilate, have a continuous potential
which causes the annihilation. They have used the same background phase
function to accelerate the annihilation process dipole4 . Recently, the tight
focusing properties of a pair of vortices have been investigated theoretically
dipole5 ; dipole6 . However, it dealt only with isopolar and dipole vortices
of first order. Here, we present a theoretical analysis of propagation of an
arbitrary vortex pair passing through an astigmatic optical system and verify
the results with experiments. We suggest that the vortices being generic to
all the waves berry , the present study can be useful to acoustic thomas and
matter waves angom also.
Since the topological charge of a vortex determines its OAM, an accurate
measurement of the topological charge is an essential and important task.
There are a number of methods to determine the charge of an optical vortex and
its sign inter ; inter1 ; inter2 ; pravin1 ; shashi ; hick ; mourka ; pravin2
. For a multi-singular beam, however, there is no method to measure the net
charge which determines the torque imparted by the field. We show that the
intensity distribution of a multi-singular beam at a predicted position beyond
a tilted lens can provide information about the net charge present in it.
## II Theory
Consider a pair of optical vortices embedded in a Gaussian beam, one with
topological charge $\epsilon_{1}m\,(\epsilon_{1}=\pm 1)$ located at
$x_{1}=-x_{0}$, $y_{1}=0$ and another with topological charge
$\epsilon_{2}n\,(\epsilon_{2}=\pm 1)$ at $x_{1}=x_{0}$, $y_{1}=0$. The complex
field distribution of the vortex pair at the waist plane of the host Gaussian
beam, with waist size $w_{0}$, is given by
$\begin{split}E_{1}(x_{1},y_{1})=(x_{1}+x_{0}+&i\epsilon_{1}y_{1})^{m}(x_{1}-x_{0}+i\epsilon_{2}y_{1})^{n}\\\
&\times\exp\left[-\left(\frac{x_{1}^{2}+y_{1}^{2}}{{w_{0}^{2}}}\right)\right].\end{split}$
(1)
The tilted lens is placed at a distance $z_{0}$ from the waist plane. The
vortex passes through the lens and travels a further distance $z$. The overall
ray transfer matrix ${\bf M_{tot}}$ is given by pravin2
${\bf M_{tot}}=\left(\begin{array}[]{cc}{\bf A}&{\bf B}\\\ -{\bf C}/f&{\bf
D}\end{array}\right)$ (2)
where ${\bf A}$, ${\bf B}$, ${\bf C}$ and ${\bf D}$ are $2\times 2$ diagonal
matrices with diagonal elements given by $a_{j}$, $b_{j}$, $c_{j}$ and $d_{j}$
respectively. Explicitly,
$\displaystyle c_{1}$ $\displaystyle=$
$\displaystyle\sec\theta,\;c_{2}=\cos\theta,\;a_{j}=1-zc_{j}/f,$
$\displaystyle d_{j}$ $\displaystyle=$ $\displaystyle 1-z_{0}c_{j}/f,\qquad
b_{j}=z_{0}+zd_{j},{j=1,2}.$ (3)
Next, we define two column vectors ${\bf r_{1}}$, ${\bf r_{2}}$ so that their
transposes are given by row vectors ${\bf r_{i}}^{T}=(x_{i},y_{i}),{i=1,2}$.
The field $E_{2}(x_{2},y_{2})$ at a distance $z$ past the lens is given by the
generalized Huygens-Fresnel integral sieg :
$E_{2}(x_{2},y_{2})=\frac{i/\lambda}{|B|^{1/2}}\int\\!\\!\\!\int
dx_{1}\,dy_{1}\,E_{1}(x_{1},y_{1})e^{-(i\pi/\lambda){\bf\phi(r_{1},r_{2})}}$
(4)
where $|B|=|b_{1}b_{2}|$ is the determinant of ${\bf B}$ and
$\displaystyle{\bf\phi(r_{1},r_{2})}$ $\displaystyle=$ $\displaystyle{\bf
r_{1}}^{T}{\bf B}^{-1}{\bf A}{\bf r_{1}}+{\bf r_{2}}^{T}{\bf D}{\bf
B}^{-1}{\bf A}{\bf r_{2}}-2{\bf r_{1}}^{T}{\bf B}^{-1}{\bf r_{2}}$ (5)
$\displaystyle=$ $\displaystyle
x_{1}^{2}a_{1}/b_{1}+y_{1}^{2}a_{2}/b_{2}+x_{2}^{2}d_{1}/b_{1}+y_{2}^{2}d_{2}/b_{2}$
$\displaystyle\mbox{}-2(x_{1}x_{2}/b_{1}+y_{1}y_{2}/b_{2}).$
The integration over $x_{1}$ and $y_{1}$ are carried out by writing
$E_{1}(x_{1},y_{1})$ as
$\displaystyle E_{1}(x_{1},y_{1})$ $\displaystyle=$
$\displaystyle\lim_{\begin{subarray}{c}t\rightarrow 0\\\ t^{\prime}\rightarrow
0\end{subarray}}\left[\frac{\partial^{m}}{\partial
t^{m}}\frac{\partial^{n}}{\partial{t^{\prime}}^{n}}\exp\left\\{f\left(t,t^{\prime}\right)\right\\}\right],$
(6a) $\displaystyle f(t,t^{\prime})$ $\displaystyle=$ $\displaystyle
t(x_{1}+x_{0}+i\epsilon_{1}y_{1})+t^{\prime}(x_{1}-x_{0}+i\epsilon_{2}y_{1})$
(6b) $\displaystyle{}-\frac{x_{1}^{2}+y_{1}^{2}}{w_{0}^{2}}.$
Using the definition of Hermite polynomial and a recurrence relation
$\displaystyle H_{n}(x)$ $\displaystyle=$
$\displaystyle\frac{\partial^{n}}{\partial{t}^{n}}\exp(2xt-t^{2})|_{t=0}$ (7a)
$\displaystyle\frac{d^{j}}{dx^{j}}H_{n}(x)$ $\displaystyle=$
$\displaystyle\frac{2^{j}n!}{(n-j)!}H_{n-j}(x)$ (7b)
we finally get
$\begin{split}E_{2}(x_{2},y_{2})=&\frac{kw_{1}w_{2}(i/2)^{m+n+1}\gamma^{m+n}}{(b_{1}b_{2})^{1/2}}\\\
&\,\times\exp\left[-\left(\beta_{1}x_{2}^{2}+\beta_{2}y_{2}^{2}\right)\right]F_{m,n}(x_{2},y_{2}),\end{split}$
(8a)
$\begin{split}F_{m,n}(x_{2},y_{2})=&\sum_{j=0}^{\min(m,n)}{m\choose
j}\,{n\choose j}\Delta^{j}j!\\\
&\,H_{m-j}[f_{1}(x_{2},y_{2})]H_{n-j}[f_{2}(x_{2},y_{2})]\end{split}$ (8b)
where, $k=2\pi/\lambda$,
$\displaystyle\frac{1}{w_{j}^{2}}$ $\displaystyle=$
$\displaystyle\frac{1}{w_{0}^{2}}+i\frac{ka_{j}}{2b_{j}},$ (9a)
$\displaystyle\gamma$ $\displaystyle=$
$\displaystyle(w_{1}^{2}-w_{2}^{2})^{1/2},$ (9b) $\displaystyle\Delta$
$\displaystyle=$
$\displaystyle-2(w_{1}^{2}-w_{2}^{2}\epsilon_{1}\epsilon_{2})/\gamma^{2},$
(9c) $\displaystyle\alpha_{j}$ $\displaystyle=$
$\displaystyle\frac{kw_{j}^{2}}{2b_{j}},$ (9d) $\displaystyle\beta_{j}$
$\displaystyle=$
$\displaystyle\left(\frac{kw_{j}}{2b_{j}}\right)^{2}+i\frac{kd_{j}}{2b_{j}},$
(9e)
and
$\displaystyle\left[\begin{array}[]{c}f_{1}(x_{2},y_{2})\\\
f_{2}(x_{2},y_{2})\end{array}\right]$ $\displaystyle=$
$\displaystyle\frac{1}{\gamma}\left[\begin{array}[]{c}\alpha_{1}x_{2}+i(\epsilon_{1}\alpha_{2}y_{2}-x_{0})\\\
\alpha_{1}x_{2}+i(\epsilon_{2}\alpha_{2}y_{2}+x_{0})\end{array}\right]$ (14)
$\displaystyle=$
$\displaystyle\frac{1}{\gamma}\left[\begin{array}[]{c}\phi_{1}(x_{2},y_{2})\\\
\phi_{2}(x_{2},y_{2})\end{array}\right]$ (17)
Eqs. (8-10) form one of our main results. It generalizes previous work dipole6
on the propagation dynamics of a vortex pair through an astigmatic system in
that the topological charges $m$ and $n$ need not be the same and can have
arbitrary integer values.
Before proceeding further, we note that the above general result includes the
following special cases: (1) For $m=n$, we get the propagation dynamics of (a)
an isopolar vortex pair if $\epsilon_{1}\epsilon_{2}=1$ and (b) a vortex
dipole if $\epsilon_{1}\epsilon_{2}=-1$; (2) For $n=0$, the $j$-sum reduces to
the $j=0$ term only, and we get the propagation dynamics for an off-center
single vortex given by
$\displaystyle E_{2}(x_{2},y_{2})$ $\displaystyle=$
$\displaystyle\frac{kw_{1}w_{2}(i/2)^{m+1}}{(b_{1}b_{2})^{1/2}}\exp\left[-\left(\beta_{1}x_{2}^{2}+\beta_{2}y_{2}^{2}\right)\right]$
(18)
$\displaystyle{}\times\gamma^{m}H_{m}[(\alpha_{1}x_{2}+i\epsilon_{1}\alpha_{2}y_{2}-ix_{0})/\gamma].$
(3) Setting $x_{0}=0$ in the above result, one immediately recovers our
previous result pravin2 for a single vortex at the origin.
The sum $F_{m,n}$ can be evaluated formally as follows. We introduce the
2-variable Hermite-Kampé de Fériet polynomials $H_{n}(x,y)$ as khan
$H_{n}(x,y)=n!\sum_{r=0}^{[n/2]}\frac{x^{n-2r}y^{r}}{(n-2r)!r!}$ (19)
in terms of which the classical Hermite polynomials $H_{n}(x)$ are given by
$H_{n}(x)=H_{n}(2x,-1).$ (20)
Next, we consider the 4-variable 2-index 1-parameter Hermite polynomials
$H_{m,n}(x,z;y,w|\tau)$ defined as khan ; khan3
$\begin{split}H_{m,n}(x,z;y,w|\tau)=&\sum_{s=0}^{\min(m,n)}\tau^{s}s!{m\choose
s}\,{n\choose s}\,\\\ &H_{m-s}(x,z)H_{n-s}(y,w).\end{split}$ (21)
It is then easy to show that
$F_{m,n}=H_{m,n}(2f_{1},-1;2f_{2},-1|\Delta)$ (22)
which has the following generating function
$\begin{split}\exp[&-(u^{2}+v^{2})+2(f_{1}u+f_{2}v)+\Delta uv]\\\
&=\sum_{m,n=0}^{\infty}\frac{u^{m}v^{n}}{m!n!}H_{m,n}(2f_{1},-1;2f_{2},-1|\Delta)\end{split}$
(23)
### II.1 Determination of net topological charge
As noted earlier pravin2 , the modulations due to the Hermite polynomial
become most prominent when $w_{2}=w_{1}^{*}$. This happens at a certain value
$z=z_{c}$. To determine $z_{c}$ and also the distance $z_{0}$ between the
waist plane and the lens, we impose the following conditions:
$\frac{ka_{1}}{2b_{1}}|_{z=z_{c}}=-\frac{ka_{2}}{2b_{2}}|_{z=z_{c}}=\frac{1}{w_{0}^{2}}$
(24)
Solving Eqs. (17) and introducing the Rayleigh range $z_{R}=kw_{0}^{2}/2$, we
get
$\displaystyle z_{0}$ $\displaystyle=$ $\displaystyle
z_{R}\left(1+\frac{2f\cos\theta}{z_{R}\sin^{2}\theta}\right)^{1/2}$
$\displaystyle z_{c}$ $\displaystyle=$
$\displaystyle\frac{z_{R}(1+\cos^{2}\theta)+z_{0}\sin^{2}\theta}{2(z_{R}/f)\cos\theta-\sin^{2}\theta}$
(25)
The first equality in Eqs. (17) ensures that $w_{2}=w_{1}^{*}$ at $z=z_{c}$
(see Eq. 9a) whereas the last equality makes many expressions appearing in
Eqs. (8-10) considerably simpler at $z=z_{c}$. Thus, at $z=z_{c}$,
$\displaystyle\Delta$ $\displaystyle=$
$\displaystyle\left\\{\begin{array}[]{ll}-2&\mbox{if
$\epsilon_{1}\epsilon_{2}=1$},\\\ -2i&\mbox{if
$\epsilon_{1}\epsilon_{2}=-1$};\end{array}\right.$ (28)
$\displaystyle\left(\begin{array}[]{c}w_{1}^{2}\\\
w_{2}^{2}\end{array}\right)$ $\displaystyle=$
$\displaystyle\frac{w_{0}^{2}}{\sqrt{2}}\left(\begin{array}[]{c}\exp(-i\pi/4)\\\
\exp(i\pi/4)\end{array}\right);$ (33) $\displaystyle\gamma$ $\displaystyle=$
$\displaystyle w_{0}\exp(-i\pi/4);$ $\displaystyle f_{1}$ $\displaystyle=$
$\displaystyle\delta_{1}x_{2}-\epsilon_{1}\delta_{2}y_{2}+(x_{0}/w_{0})\exp(-i\pi/4)$
$\displaystyle f_{2}$ $\displaystyle=$
$\displaystyle\delta_{1}x_{2}-\epsilon_{2}\delta_{2}y_{2}-(x_{0}/w_{0})\exp(-i\pi/4)$
(34)
where
$\delta_{j}=\frac{kw_{0}}{2\sqrt{2}b_{i}}$ (35)
#### II.1.1 Vortices with topological charges of the same sign
Suppose $\epsilon_{1}=\epsilon_{2}=1$. Then, $f_{1}=\theta_{-}+\theta_{0}$ and
$f_{2}=\theta_{-}-\theta_{0}$ where,
$\displaystyle\theta_{-}$ $\displaystyle=$
$\displaystyle\delta_{1}x_{2}-\delta_{2}y_{2}$ $\displaystyle\theta_{0}$
$\displaystyle=$ $\displaystyle(x_{0}/w_{0})\exp(-i\pi/4).$ (36)
Note that the dependence on $x_{2}$ and $y_{2}$ is in the form $\theta_{-}$
only.
For a small separation between the vortices, one can expand the Hermite
polynomials appearing in Eq.(8) as functions of $x_{0}/w_{0}$ by using the
formula
$H_{n}(x+y)=H_{n}(x)+2nyH_{n-1}(x)+O(y^{2}).$ (37)
Substituting in (8) and using the summation rule magnus
$\sum_{r=0}^{\min(m,n)}(-2)^{r}r!{m\choose r}\,{n\choose
r}\,H_{m-r}(x)H_{n-r}(x)=H_{m+n}(x),$ (38)
we get
$F_{m,n}=H_{m+n}(\theta_{-})+2\theta_{0}(m-n)H_{m+n-1}(\theta_{-})+O(\theta_{0}^{2}).$
(39)
For $\epsilon_{1}=\epsilon_{2}=-1$, $\theta_{-}$ will change to
$\theta_{+}=\delta_{1}x_{2}+\delta_{2}y_{2}$ in the above expressions.
#### II.1.2 Vortices with topological charges of opposite signs
Suppose $\epsilon_{1}=1$, $\epsilon_{2}=-1$. In this case,
$\displaystyle f_{1}$ $\displaystyle=$ $\displaystyle\theta_{-}+\theta_{0}$
$\displaystyle f_{2}$ $\displaystyle=$ $\displaystyle\theta_{+}-\theta_{0}.$
(40)
Note that in this case, the dependence on $x_{2}$ and $y_{2}$ is in the form
$\theta_{\pm}=\delta_{1}x_{2}\pm\delta_{2}y_{2}$.
For a small separation between the vortices, we can proceed as in the previous
section, to get
$\displaystyle F_{m,n}=H_{m,n}(2\theta_{-},-1;2\theta_{+},-1|-2i)$ (41)
$\displaystyle{}+2\theta_{0}H_{m-1,n}(2\theta_{-},-1;2\theta_{+},-1|-2i)$
$\displaystyle{}-2\theta_{0}H_{m,n-1}(2\theta_{-},-1;2\theta_{+},-1|-2i)$
$\displaystyle{}+O(\theta_{0}^{2}).$
### II.2 Propagation dynamics away from $z=z_{c}$
As $|z-z_{c}|$ increases, the absolute value of $|\gamma|$ falls off rapidly
and the modulations due to the Hermite polynomials fade away quickly. Using
the limiting form $\lim_{\gamma\rightarrow 0}H_{m}(x/\gamma)=(2x/\gamma)^{m}$,
we can write $E_{2}(x_{2},y_{2})$ in terms of incomplete two-variable Hermite
polynomials $h_{m,n}(x,y|\tau)$, which are defined as khan ; khan4
$\displaystyle
h_{m,n}(x,y|\tau)=m!n!\sum_{j=0}^{\min(m,n)}\frac{\tau^{j}x^{m-j}y^{n-j}}{j!(m-j)!(n-j)!}$
(44) $\displaystyle=$
$\displaystyle\left\\{\begin{array}[]{ll}m!\tau^{m}x^{n-m}L_{m}^{(n-m)}(-xy/\tau),\qquad
n>m,\\\ n!\tau^{n}y^{m-n}L_{n}^{(m-n)}(-xy/\tau),\qquad
m>n.\end{array}\right.$
Thus $E_{2}(x_{2},y_{2})$ reduces to
$\displaystyle
E_{2}(x_{2},y_{2})=\frac{kw_{1}w_{2}i^{m+n+1}}{2(b_{1}b_{2})^{1/2}}\exp\left[-\left(\beta_{1}x_{2}^{2}+\beta_{2}y_{2}^{2}\right)\right]$
$\displaystyle\times\left\\{\begin{array}[]{ll}m!\tau^{m}\phi_{1}^{n-m}L_{m}^{(n-m)}(-\phi_{1}\phi_{2}/\tau),\qquad
n>m,\\\ n!\tau^{n}\phi_{2}^{m-n}L_{n}^{(m-n)}(-\phi_{1}\phi_{2}/\tau),\qquad
m>n.\end{array}\right.$ (47)
where $\phi_{j}$ are as in Eq. (10) and
$\tau=-(w_{1}^{2}-w_{2}^{2}\epsilon_{1}\epsilon_{2})/2$. In what follows, we
will experimentally demonstrate the validity of our theoretical results.
## III Experiment
The experimental set up is shown in Fig. 1. Suitable phase masks for creating
vortex pairs are produced by using computer generated holography (CGH)
technique cgh and sent to a spatial light modulator (SLM) via a computer. The
SLM is illuminated by an intensity stabilized He-Ne laser (Spectra-Physics,
Model 117A) of power $1$ mW and wavelength $632.8$ nm to produce the desired
vortex pair. The vortex pair is selected with an aperture (A) and passed
through a spherical bi-convex lens of focal length $50$ cm which is tilted by
an angle $6^{\circ}$. The tilting of the lens has been done with a rotational
stage with least count of $0.1^{\circ}$. The aperture is at a distance
$z_{1}=90$ cm in front of the SLM. We use the method described in sirohi , to
find that the Gaussian laser beam hosting the selected vortex pair has a beam
waist $0.186$ mm at a virtual point which is at a distance of $z_{2}=60.8$ cm
behind the SLM. The distance between the lens and the aperture is $z_{3}=245$
cm. Thus the total distance traveled by the vortex pair from the waist plane
to the lens is $z_{0}=z_{1}+z_{2}+z_{3}=395.8$ cm. The resultant intensity
patterns are recorded by a CCD camera (MediaCybernetics, Evolution VF cooled
Color Camera) placed at a distance $z$ past the lens.
Figure 1: (Colour online) Experimental setup for the determination of the net
charge of an arbitrary vortex pair embedded in a Gaussian beam
## IV Intensity pattern at $z=z_{c}$ and determination of net topological
charge
In this section, we determine the net topological charge of the vortex pair
from its intensity distribution at $z=z_{c}$. The predicted value of $z_{c}$
from Eq. (25) is $57.2$ cm which is close to the experimentally observed value
of $56.3$ cm. In the intensity patterns, with reference to Eq. (LABEL:vor),
the vortex on the left ($x_{1}=-x_{0}$) has a charge $\epsilon_{1}m$ and the
vortex on the right ($x_{1}=x_{0}$) has a charge $\epsilon_{2}n$. The
corresponding figure is labelled as ($\epsilon_{1}m,\,\epsilon_{2}n$).
### IV.1 Vortices with topological charges of the same sign
Figure 2: The theoretical (first two rows) and experimental (last two rows)
results for the intensity patterns of a vortex pair with topological charges
of the same sign, at $z=z_{c}$ for $x_{0}=0.1w_{0}$.
Fig. 2 shows the theoretical (first two rows) and experimental (last two rows)
images for the intensity patterns of a pair of vortices with the same sign
($\epsilon_{1}=\epsilon_{2}=1$) but different magnitudes $m$ and $n$ with the
separation parameter set at $x_{0}=0.1w_{0}$.
For small separation $x_{0}$, these patterns can be explained by Eq. (24).
Since the first term in (24) is the leading term, one can obtain the net
charge $m+n$ by noting that there are $m+n+1$ bright stripes in the intensity
distribution. These stripes are parallel to one another and lie along a line
that is neither horizontal nor vertical, but tilted in a clockwise direction
almost along a diagonal as the dependence on $x$ and $y$ is through a single
variable $\theta_{-}=\delta_{1}x_{2}-\delta_{2}y_{2}$ and
$\delta_{1}\sim\delta_{2}$. However, interference with the second term will
lead to a slightly asymmetric distribution of brightness among the stripes. As
is clear from the second term in (24), this asymmetry depends on the
difference between the magnitude of charges and the separation between them.
Additionally, when the vortices swap their positions as in (4,1) and (1,4),
the lower half of the pattern becomes the mirror image of the upper half and
vice-versa. For $m=n$ as in (4,4), the two halves have identical intensity
patterns. If the charge of each vortex were negative
($\epsilon_{1}=\epsilon_{2}=-1$), then $\theta_{-}$ would be replaced by
$\theta_{+}=\delta_{1}x_{2}+\delta_{2}y_{2}$ and the bright stripes would be
tilted in an anti-clockwise fashion (not shown).
Figure 3: The theoretical (first two rows) and experimental (last two rows)
results for the intensity patterns of an off-axis vortex of charge 2, at
$z=z_{c}$ for different values of $x_{0}$ as labelled in the figures.
To investigate the effect of the separation parameter $x_{0}$, we have also
studied the propagation of an off-axis vortex of charge 2 through the tilted
lens. The corresponding theoretical (first two rows) and experimental (last
two rows) results for the intensity patterns at $z=z_{c}$ are shown in Fig. 3.
From the images, it is clear that the intensity of one of the outer lobes
increases as the vortex moves farther away from the center and the remaining
lobes lose their intensity. In the notation of this section, this off-axis
vortex can be labeled as (2,0) for $x_{0}=|x_{0}|$. Consequently, when $x_{0}$
becomes negative, the vortex is identically described as (0,2) with
$x_{0}=|x_{0}|$ and the pattern flips diagonally. The situation is analogous
to the case of (4,1) and (1,4) as described in the previous paragraph
### IV.2 Vortices with topological charges of opposite signs
Fig. 4 shows the theoretical (first two rows) and experimental (last two rows)
images corresponding to opposite singularities ($\epsilon_{1}\epsilon_{2}=-1$)
for separation parameter $x_{0}=0.1w_{0}$ and topological charges as shown in
the images.
Figure 4: The theoretical (first two rows) and experimental (last two rows)
results for the intensity patterns of a vortex pair with topological charges
of opposite signs, at $z=z_{c}$ for $x_{0}=0.1w_{0}$.
For small values of $m$ and $n$, these patterns can be explained by expanding
the Hermite polynomials in power series. The calculation would be long and
tedious. Instead, we make the following empirical observation. If $m\neq n$,
the pattern has a rectangular ‘razor-blade’ structure which is tilted
clockwise (anti-clockwise) if the net charge is positive (negative). On closer
observation, we note that there are $m$ bright spots on two parallel sides and
$n$ bright spots on the remaining two parallel sides. Thus, for vortex dipoles
($m=n$) the patten is square with its corners in the east, west, north and
south directions, each side having $m=n$ bright spots.
As far as we know, Fig. 4 represents the first optical realization of the
4-variable 2-index 1-parameter Hermite polynomials $H_{m,n}(x,z;y,w|\tau)$
modulated by an elliptical Gaussian beam (see Eqs. 8, 15 and 26).
Figure 5: The theoretical (top row) and experimental (bottom row) results for
the intensity patterns of a dipole vortex of charge $(2,-2)$, at $z=z_{c}$ for
different values of $x_{0}$ as labelled in the figures.
In Fig. 5, we show the evolution of a dipole vortex of charge $(2,-2)$ as a
function of separation between the two vortices. For small separation, the
intensity distribution is symmetric in both the transverse directions. As the
separation is increased, the pattern becomes asymmetric. When the separation
parameter $x_{0}$ becomes negative, the vortex is identically described as
$(-2,2)$ with $x_{0}=|x_{0}|$ and the pattern flips vertically.
Figure 6: Theoretical intensity patterns of a vortex pair of different
charges (as given on the top) at various values of the propagation distance
$z$ ( as given on the left).
## V Propagation dynamics away from $z=z_{c}$
As we move away from the point $z=z_{c}$, the modulations due to the Hermite
polynomials disappear quickly. The propagation dynamics is now governed by Eq.
(28). The theoretical and corresponding experimental intensity patterns for
various values of $z$ are shown in Figs. 6 and 7 respectively. The intensity
patterns are, in general, elliptical. Far away from $z_{c}$, all patterns
become circularly symmetric as $\alpha_{1}\to\alpha_{2}$ and
$\beta_{1}\to\beta_{2}$ pravin2 .
Figure 7: The experimental images corresponding to Fig. 6.
## VI Conclusions
We have studied the propagation of a bi-singular beam with arbitrary
topological charges through a tilted lens and used it to find the net
topological charge of the beam. This may provide information about the net
torque generated by the optical field. This method can be realized easily in
the laboratory as it needs just a single tilted lens. Vortices being generic
to all the waves, this study may be useful for other systems like acoustic and
matter waves.
## References
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|
arxiv-papers
| 2013-09-12T10:23:54 |
2024-09-04T02:49:50.815363
|
{
"license": "Public Domain",
"authors": "Salla Gangi Reddy, Shashi Prabhakar, A Aadhi, J. Banerji, and R. P.\n Singh",
"submitter": "Salla Gangi Reddy",
"url": "https://arxiv.org/abs/1309.3095"
}
|
1309.3163
|
11institutetext: Department of Physics and Centre for Computational Science
and Engineering, National University of Singapore, Singapore 117542, Republic
of Singapore
Department of Physics and Institute of Theoretical Physics and Astrophysics,
Xiamen University, Xiamen 361005, China
NUS Graduate School for Integrative Sciences and Engineering, Singapore
117456, Republic of Singapore
NUS-Tongji Center for Phononics and Thermal Energy Science and Department of
Physics, Tongji University, 200092 Shanghai, PR China
Heat conduction Solid surfaces and solid-solid interfaces: structure and
energetics Scattering by phonons, magnons, and other nonlocalized excitations
# Nonlinearity enhanced interfacial thermal conductance and rectification
Lifa Zhang 11 Juzar Thingna 11 Dahai He 22 Jian-Sheng Wang 11 Baowen Li
11334411223344
###### Abstract
We study the nonlinear interfacial thermal transport across atomic junctions
by the quantum self-consistent mean field (QSCMF) theory based on
nonequilibrium Green’s function approach; the QSCMF theory we propose is very
precise and matches well with the exact results from quantum master equations.
The nonlinearity at the interface is studied by effective temperature
dependent interfacial coupling calculated from the QSCMF theory. We find that
nonlinearity can provide an extra channel for phonon transport in addition to
the phonon scattering which usually blocks heat transfer. For weak linearly
coupled interface, the nonlinearity can enhance the interfacial thermal
transport; with increasing nonlinearity or temperature, the thermal
conductance shows nonmonotonical behavior. The interfacial nonlinearity also
induces thermal rectification, which depends on the mismatch of the two leads
and also the interfacial linear coupling.
###### pacs:
44.10.+i
###### pacs:
68.35.-p
###### pacs:
72.10.Di
## 1 Introduction
In modern electronics, due to the rapid increasing power density, accumulation
of heat becomes an obstacle for further progress of microelectronic devices;
thus the heat dissipation and manipulation has been recognized to be a crucial
issue in information and energy technologies [1]. Especially, as the
dimensions of materials shrink into the nanoscale, interfaces dramatically
affect the thermal transport [2, 3, 4, 5] making it a lucrative field to
explore. At a rough interface the atomic mixing can enhance the thermal
transport [6]; however, the behavior of the nonlinearity at the interface is
not clear, and people do not know whether the nonlinearity can enhance the
phonon transport. Recent progress in functional thermal devices [7, 8, 9, 10],
makes the emerging new field – phononics very attractive [11]. In phononics,
the most fundamental property of phononic devices is thermal rectification,
which is known to be realized by combining the system inherent anharmonicity
with structural asymmetry [12, 13]. Whether the interface itself can induce
thermal rectification is still an open question; if yes, the property of the
interfacial rectification is quite interesting and helpful for both theorists
and experimentalists.
To investigate the thermal transport across interface, the most widely applied
models are the acoustic mismatch model [14] and the diffuse mismatch model
[15]. Both models offer limited accuracy in nanoscale interfacial resistance
predictions [16, 17] because they make simple assumptions and neglect atomic
details of actual interfaces. Classical molecular dynamics simulation is
another widely used method in phonon transport and has been applied to
interfacial thermal transport [18, 19, 20, 21, 22]; however due to its
classical nature, it is not accurate below the Debye temperature and can not
capture the quantum effects. To study the nonlinear (anharmonic) thermal
transport, the effective phonon theory has been recently introduced in some
dynamical models [23, 24, 25]; and the quantum correction one [26] can be used
to study the low temperature thermal transport. Despite their successes, such
theories can not be well applied to nonlinear interfacial transport due to the
inherent weak system-bath coupling assumption required for the validity of
Feynman-Jensen inequality [27]. Another effective approach, the nonequilibrium
Green’s function method which originates from the study of electronic
transport [28], has been applied to study the quantum phonon transport [30,
31, 32], phonon Hall effect [33] and topological magnon insulator [34].
In this paper, based on the nonequilibrium Green’s function method, to avoid
the perturbation approximation we develop the QSCMF theory for the nonlinear
thermal transport, which can be applied to thermal transport in an arbitrary
strength nonlinear interface. Then we study the interfacial thermal transport
for a model as shown in Fig.1(a); thermal conductance and rectification across
the interface are studied with an effective temperature-dependent-harmonic
interfacial coupling calculated form the QSCMF theory.
## 2 Model and Hamiltonian
We study the interfacial thermal transport with nonlinear coupling at the
solid-solid interface as shown in Fig. 1(a). To manifest the effect of the
interface we exclude the nonlinear phonon transport in the two materials, and
there is no disorder nor defect in the whole system, thus the only thermal
resistance comes from the interface. Such model really uncovers the thermal
transport properties of the interface. To study the longitudinal transport,
that is, the cross-plane interfacial transport, we simplify the problem
further to one dimensional interface as shown in Fig. 1(b), where two linear
semi-infinite atomic chains (solid-line regimes) connect each other by the
interfacial linear coupling $k_{12}$ and nonlinear interaction $\lambda$,
which is similar to the one-dimensional interfacial model in Ref.[5, 3] where
there is only linear coupling at the interface. The Hamiltonian of the total
system is
Figure 1: (Color online) (a) Heat transport in a solid-solid interface. The
arrow shows the heat transport from the hot side to the cold side. (b) The
atomic junction model of the solid-solid interface. The solid line regions are
two semi-infinite atomic chains which are coupled by a harmonic spring with
strength $k_{12}$. In addition to which, the two regions also have a fourth
order nonlinear coupling $\lambda$. For the two semi-infinite chains, the mass
and spring constant are $m_{1}$, and $k_{1}$, $m_{2}$, and $k_{2}$,
respectively. The interface model can be partitioned to three parts, the
center (dashed line) and the leads (dash-dotted lines) with temperatures
$T_{1}$ and $T_{2}$.
$H=H_{1}+H_{2}+\frac{1}{2}k_{12}(x_{1,1}-x_{2,1})^{2}+\frac{1}{4}\lambda(x_{1,1}-x_{2,1})^{4},$
(1)
with
$H_{\alpha}=\sum\limits_{i=1}^{N_{\alpha}}{\frac{1}{2}m_{\alpha}\dot{x}_{\alpha,i}^{2}}+\sum\limits_{i=1}^{N_{\alpha}-1}{\frac{1}{2}k_{\alpha}(x_{\alpha,i}-x_{\alpha,i+1})^{2}},$
(2)
where $N_{\alpha}\rightarrow\infty$, $\alpha=1,2$. In our model, the
scattering for phonons only comes from the interface while the phonon
transport in the two semi-infinite leads is ballistic. Thus we can partition
the system into three parts ($L$, $C$, $R$), where the atoms at the interface
are regarded as center (dashed-line part) and leads $L$ and $R$ (dash-dotted-
line regimes) are harmonic as shown in Fig.1 (b). Using a mass-normalized
displacement $u_{j}=\sqrt{m_{j}}x_{j,1}$, the center Hamiltonian can be
written as
$H_{C}=\frac{1}{2}\dot{U}_{c}^{T}\dot{U}_{c}+\frac{1}{2}U_{c}^{T}K^{C}U_{c}+\frac{1}{4}\sum\limits_{i,j,k,l=1}^{2}{T_{ijkl}}u_{i}u_{j}u_{k}u_{l},$
(3)
where $U_{c}=(u_{1},u_{2})^{T}$, $T_{ijkl}=(-1)^{i+j+k+l}\lambda$ and
$K^{C}=\left({\begin{array}[]{*{20}c}\frac{k_{1}+k_{12}}{m_{1}}&\frac{-k_{12}}{\sqrt{m_{1}m_{2}}}\\\
\frac{-k_{12}}{\sqrt{m_{1}m_{2}}}&\frac{k_{12}+k_{2}}{m_{1}}\\\
\end{array}}\right).$ (4)
For the leads, the Hamiltonian is written as
$H_{\alpha}=\frac{1}{2}\dot{U}_{\alpha}^{T}\dot{U}_{\alpha}+\frac{1}{2}U_{\alpha}^{T}K^{\alpha}U_{\alpha}$
with its coupling to the center $H_{\alpha C}=U_{\alpha}^{T}V^{\alpha
C}U_{C}$, $\alpha=L,R$. Here the center-lead coupling is the same as the
inter-atomic spring constant in the corresponding bath.
## 3 Quantum Self-Consistent Mean Field Theory
We discuss the QSCMF based on NEGF method for a general system where the
center hamiltonian has a fourth-order nonlinear interaction as given in Eq.
(3). The equation of motion of Green’s function [32], without the
nonlinearity, is
$(\frac{\partial^{2}}{{\partial\tau^{2}}}+K^{C})G_{0}(\tau\tau^{\prime})=-I\delta(\tau-\tau^{\prime})-\int
d\tau^{\prime\prime}\Sigma(\tau\tau^{\prime\prime})G_{0}(\tau^{\prime\prime}\tau^{\prime}),$
where $\Sigma(\tau\tau^{\prime\prime})$ is the self energy due to the center-
lead coupling. With the nonlinearity, the full Green’s function has the
equation of motion as
$\begin{split}\frac{{\partial^{2}}}{{\partial\tau^{2}}}G_{im}(\tau\tau^{\prime})+\sum\limits_{j}{K_{ij}^{C}}G_{jm}(\tau\tau^{\prime})+\sum\limits_{jkl}{T_{ijkl}G_{jklm}(\tau\tau\tau\tau^{\prime})}\\\
=-\delta(\tau-\tau^{\prime})\delta_{im}-\sum\limits_{j}{\int
d\tau^{\prime\prime}\Sigma_{ij}(\tau\tau^{\prime\prime})G_{jm}(\tau^{\prime\prime}\tau^{\prime})}\end{split}$
(5)
where $G_{jklm}(\tau_{j}\tau_{k}\tau_{l}\tau_{m})$ is the four-point Green’s
function. Under the self-consistent mean field approximation, the four-point
Green’s function can be represented by the two-point Green’s function, and we
have
$-i\frac{1}{\hbar}G(\tau_{j}\tau_{k}\tau_{l}\tau_{m})=G(\tau_{j}\tau_{k})G(\tau_{l}\tau_{m})+G(\tau_{j}\tau_{l})G(\tau_{k}\tau_{m})+G(\tau_{j}\tau_{m})G(\tau_{k}\tau_{l}).$
Thus the full Green’s function satisfies
$\begin{split}\frac{{\partial^{2}}}{{\partial\tau^{2}}}G_{im}(\tau\tau^{\prime})+\sum\limits_{j}{K_{ij}^{C}}G_{jm}(\tau\tau^{\prime})\\\
+3i\hbar\sum\limits_{jkl}{T_{ijkl}G_{kl}(0)G_{jm}(\tau\tau^{\prime})}+\sum\limits_{j}{\int
d\tau^{\prime\prime}{\Sigma_{ij}(\tau\tau^{\prime\prime})G_{jm}(\tau^{\prime\prime}\tau^{\prime})}}\\\
=-\delta(\tau-\tau^{\prime})\delta_{im},\end{split}$ (6)
by using the symmetry of $T_{ijkl}$ with respect to the permutation of the
indices. Thus we introduce an effective dynamic matrix
$\tilde{K}_{ij}^{C}=K_{ij}^{C}+3i\hbar\sum\limits_{kl}{T_{ijkl}G_{kl}(0)}=K_{ij}^{C}+3\sum\limits_{kl}{T_{ijkl}\left\langle{u_{k}u_{l}}\right\rangle}.$
(7)
The nonlinearity only has the effect to modulate the dynamic matrix. It is
important to note that the results are independent of the partition size of
the center due to $T_{ijkl}=0$ every where else except the two atoms at the
interface.
Equation (7) together with
$G^{r}=[(\omega+i\eta)^{2}-\tilde{K}_{ij}^{C}-\Sigma^{r}]^{-1},\;G^{<}=G^{r}\Sigma^{<}G^{a},\;\left\langle{u_{k}u_{l}}\right\rangle=i\hbar
G_{kl}^{<}(t=0)=2i\hbar\int_{0}^{\infty}{G_{kl}^{<}(\omega)d\omega/(2\pi)}$
can be self-consistently calculated. Since the problem is now effectively
harmonic, the heat current still satisfies the Landauer formula
$J=\int_{0}^{\infty}{\frac{{d\omega}}{{2\pi}}\hbar\omega
T[\omega](f_{L}-f_{R})}$,
$f_{\alpha}=1/(e^{\hbar\omega/(k_{B}T_{\alpha})}-1)$, and thermal conductance
is defined as $\sigma=|J/(T_{L}-T_{R})|$, while the transmission
$T[\omega]={\rm Tr}(G^{r}\Gamma_{L}G^{a}\Gamma_{R})$ is temperature dependent.
In later calculation, we will set $\hbar=1,k_{B}=1$ for simplicity. For the
conversion from the dimensionless unit to physical units, we take energy unit
$[E]=1$ mev, length unit is $[L]=1{\rm\AA}$, then temperature unit $[T]$ is
11.6 K, thermal conductance $[\sigma]=20.9{\rm nW/mK}$, the spring constant
unit $[k]=1$ mev/${\rm\AA}^{2}$ and nonlinearity unit $[\lambda]=1$
mev/${\rm\AA}^{4}$.
Applying our QSCMF theory to the nonlinear interface problem of Eq. (1), we
find that the nonlinearity plays a role to modulate the interfacial linear
coupling $k_{12}$, the effective one is
$k_{12\rm{eff}}=k_{12}+3\lambda\bigl{(}\frac{{\left\langle{u_{1}^{2}}\right\rangle}}{{m_{1}}}-2\frac{{\left\langle{u_{1}u_{2}}\right\rangle}}{{\sqrt{m_{i}m_{j}}}}+\frac{{\left\langle{u_{2}^{2}}\right\rangle}}{{m_{2}}}\bigr{)}.$
(8)
Thus it is clear that all the scattering occurs only at the interface, since
all other parts are harmonic.
Figure 2: (Color online) Comparison of the results for thermal conductance of
the interface as illustrated in Fig. 1(b) between QSCMF theory and quantum
master equation method. We add a small onsite potential $0.01$ to each atom in
both leads. $m_{1}=m_{2}=1.0$, $k_{1}=k_{2}=1.0$, and
$K_{11}^{C}=1.5,K_{22}^{C}=1.5$, $K_{12}^{C}=-0.5,K_{21}^{C}=-0.5$. The
system-bath coupling is $0.1$ .
## 4 Comparison with the Quantum Master Equation
While the effective phonon theory and its quantum correction one are valid in
the linear response region at weak system-bath coupling, our QSCMF theory can
study the nonequilibrium thermal transport under larger temperature bias at
any system-bath coupling. For anharmonic systems with arbitrary strength of
anharmonicity under the weak system-bath coupling approximation, Redfield
quantum master equation is very suited solution to study the thermal transport
[12, 35]. In our interface problem, we partition the atoms in the interface as
a center, and the center-lead coupling is the same as the inter-atomic spring
constant in the corresponding bath. Thus for such system, the system-bath
coupling is always strong such that the quantum master equation cannot be
applied which is limited in the weak system-bath coupling limit. However, the
QSCMF method can be applied.
We numerically compare our theory with the quantum master equation at a weak
system-bath coupling, as shown in Fig. 2. In Fig. 2, the system-bath coupling
is $0.1$ while the inter-atomic spring constant in the leads are $1.0$ such
that the weak-system-bath-coupling condition is well satisfied. A small onsite
potential $0.01$ is added to each atom in both leads to avoid the divergency
in quantum master equation method. From Fig. 2 we find that the results from
QSCMF perfectly matches those from the quantum master equation method for
different nonlinearity. Therefore the QSCMF method is verified by the quantum
master equation method at weak system-bath coupling; however, the QSCMF is not
limited in this weak coupling, which can be applied to arbitrary-system-bath-
coupling systems.
## 5 Numerical Results on Interface thermal Transport
### 5.1 Nonlinearity Suppressed Thermal Transport in Homogenous Systems
Using the QSCMF theory we proposed above, the interface nonlinearity can be
studied for the interfacial thermal transport. For the homogeneous lattice
with $k_{1}=k_{2}=k_{12}$, we calculate the interfacial thermal conductance
for different nonlinearity $\lambda$ as shown in Fig. 3(a). With zero
nonlinearity, the thermal conductance increases with the temperature
increasing, and tends to a constant due to the saturate of phonon modes
contributing to the thermal transport. However, with nonzero interfacial
nonlinearity the thermal conductance decreases at high temperatures due to the
dominant scattering coming from the nonlinear interface coupling. In the low-
temperature regime, the thermal conductance almost coincides with the
ballistic transport and the nonlinearity almost has no effect on thermal
transport. With increasing temperature, more phonon modes are excited which
wins the suppressing effect nonlinear scattering such that the conductance
increases. At certain temperature, the conductance arrives its maximum, after
which it will decrease since the nonlinear scattering effect which defeats the
enhancement effect from more excited phonon modes. With increasing
nonlinearity the thermal conductance decreases due to the larger phonon
scattering at the interface. As shown in Fig. 3(b), the nonlinearity always
decreases the thermal transport for the homogeneous systems. The larger
interfacial nonlinearity makes the system more nonhomogeneous such to induce
more scattering to the phonon transport.
Figure 3: (Color online) Interfacial thermal conductance as a function of $T$
($T=(T_{L}+T_{R})/2$)(a) and $\lambda$ (b). $m_{1}=m_{2}=1.0$,
$k_{1}=k_{2}=1.0$ and $k_{12}=1.0$. In (b), the solid line ($T=0.1$)
corresponds to left scale, and the dashed line ($T=0.5$) and dotted line
($T=2.5$) correspond to right scale.
### 5.2 Nonlinearity Enhanced Interfacial Thermal Conductance
In the weak interfacial coupling regime, that is, $k_{12}<k_{1}=k_{2}$, we
find that the nonlinear interaction at interface can enhance the thermal
transport as shown in Fig 4. With increasing nonlinearity, the interfacial
conductance increases first; after certain maximum the conductance will
decrease, as shown in Fig 4(a). The maximum conductance coincides with the
point where the effective coupling equals to $1$ and the whole system is
homogeneous. If $\tilde{k}_{12}$ increases further the conductance decreases
due to the larger scattering at interface. At a fixed linear interfacial
coupling $k_{12}<1$, the interfacial nonlinearity makes the effective one
$\tilde{k}_{12}$ larger than $k_{12}$. A larger $\tilde{k}_{12}$ reduces the
difference between the interface and the leads, thus decreases the phonon
scattering and allows more phonons to transmit through the interface.
Therefore the nonlinearity introduces an extra channel to transport phonons,
which enhances the thermal transport. In Fig 4 (b), the maximum of thermal
conductance does not coincide with the place of $\tilde{k}_{12}=1$, this is
mainly because that the increase in temperature causes more phonon modes to
transport so as to delay the maximum of conductance. For a larger
nonlinearity, the maximum of thermal conductance is delayed further.
Figure 4: (Color online) Interfacial thermal conductance and effective
interfacial coupling as functions of $\lambda$ (a) and $T$ (b).
$m_{1}=m_{2}=1.0$, $k_{1}=k_{2}=1.0$ and $k_{12}=0.80$. For both (a) and (b),
the solid, dashed and dotted lines correspond to left scale (interfacial
thermal conductance), and the scatter points (square, circle and triangle
symbols) correspond to right scale (effective interfacial coupling). The lines
with two arrows are drawn to guide the eyes.
For a hetero-junction, that is, the two mismatched leads, the thermal
conductance can also be enhanced by interfacial nonlinearity. We know that the
thermal conductance reaches a maximum when the interfacial coupling equals the
harmonic average of the spring constants of the two semi-infinite chains [5].
As shown in Fig. 5(a) the conductance of two general hetero-junctions gets to
the maximum when the effective coupling $k_{12\rm{eff}}$ equals to
$k_{12m}=2k_{1}k_{2}/(k_{1}+k_{2})$; after which the thermal conductance
decreases. The interfacial conductance can be enhanced in a wide nonlinearity
range, as shown in Fig. 5(a). The enhancement of the conductance due to the
nonlinearity depends on the interfacial linear coupling. In order to show
this, we plot the thermal conductance as a function of the linear coupling
$k_{12}$ in Fig. 5(b). For a fixed nonlinearity, the thermal conductance
enhances largest at zero $k_{12}$, where the only channel to transport phonon
comes from the nonlinearity. With increasing $k_{12}$ the conductance
increases, but the effective coupling due to the nonlinearity
$k_{12\rm{eff}}-k_{12}$ decreases. This effect causes the overall enhancement
to decrease. The inset of Fig. 5(b) shows that the nonlinear interfacial
conductance equals the linear thermal case just before $k_{12}=k_{12m}=1.5$,
where the scattering from nonlinearity cancels the contribution from the extra
channel introduced by the nonlinearity. At $k_{12}=k_{12m}=1.5$, the linear
conductance maximizes and the nonlinearity only adds to the phonon scattering,
thus the conductance of nonlinear interface is small. If $k_{12}>>\lambda$,
the effective coupling from the nonlinearity $k_{12\rm{eff}}-k_{12}$ nearly
vanishes making the interfacial nonlinearity not important anymore, thus the
conductance is almost the same with the ballistic one.
Figure 5: (Color online) (a) Thermal conductance (left scale) and effective
interfacial coupling (right scale) as functions of $\lambda$.
$k_{1}=m_{1}=1.0$, $k_{12}=0.8$, the dashed and dotted lines correspond to
interfacial thermal conductance, and the scatter pionts (circle or triangle
symbols) correspond to effective interfacial coupling; the lines with two
arrows are drawn to guide the eyes. (b) The effective coupling from the
nonlinearity $k_{12\rm{eff}}-k_{12}$ and interfacial thermal conductance
$\sigma$ as functions of interfacial linear coupling. The dashed line
corresponds to the left scale ($k_{12\rm{eff}}-k_{12}$) and solid and dotted
lines correspond to right scale ($\sigma$). $k_{1}=m_{1}=1.0$,
$k_{2}=3.0,m_{2}=2.0$, $T_{0}=0.5$; the inset is the zoom-in plot around
$k_{12}=1.5$
### 5.3 Nonlinearity Induced Interfacial Thermal Rectification
From our QSCMF approach, the effective interfacial coupling is temperature
dependent; if we reverse the two temperatures of the leads it would be
different so that we can observe thermal rectification. If the interface
coupling is linear, then Landauer formula applies and reverse temperature only
changes the sign of the heat current, thus there is no rectification. The
rectification is defined as [36]:
$R=(J_{+}-J_{-})/{\rm max}\\{J_{+},J_{-}\\},$ (9)
where $J_{+}$ is the-forward direction heat flux when
$T_{L}=T_{h},T_{R}=T_{c}$, and $J_{-}$ is that of the backward direction when
$T_{L}=T_{c},T_{R}=T_{h}$. Here $T_{h}$ and $T_{c}$ correspond to the
temperatures of the hot and cold baths, respectively.
With the asymmetric structure the nonlinear interface shows rectification,
which depends on the interfacial nonlinearity $\lambda$ and the linear
coupling $k_{12}$ as shown in Fig. 6. If the nonlinearity is zero, there is no
rectification. With increasing the nonlinearity, the effective coupling will
be different in the forward and backward transport which causes the
rectification to increase. When the nonlinearity increases further, the
effective coupling will monotonically increase. Therefore the scattering from
the interface will play an important role to decrease the flux difference
between the forward and backward ones, which causes the rectification to
decrease as shown in Fig. 6(a). With larger linear interfacial coupling, the
contribution to the phonon transport from the channel of $k_{12}$ is larger as
compared to the extra channel provided by the nonlinearity $\lambda$.
Therefore the relative nonlinear effect is weakened and the rectification is
reduced as shown in Fig. 6(b). In the opposite limit, if $k_{12}<<\lambda$ the
rectification almost keeps fixed, when the nonlinearity dominates the thermal
transport and the difference of the effective coupling between the forward and
backward flow has almost no changes. When $k_{12}\sim\lambda$, the
rectification decreases fast due to the fast decrease of the value of
$k_{12\rm{eff}}-k_{12}$. If $k_{12}>>\lambda$, the rectification decreases to
almost zero since the effective coupling from the nonlinearity
$k_{12\rm{eff}}-k_{12}$ nearly vanishes making the nonlinearity not important
anymore.
Figure 6: (Color online) (a) Interfacial thermal rectification as functions
of $\lambda$ for different $k_{12}$. (b) Interfacial thermal rectification
(left scale) and the heat current (right scale) as functions of $k_{12}$. For
both (a) and (b), $k_{1}=m_{1}=1.0$, $k_{2}=3.0,m_{2}=2.0$, and
$T_{+}=T_{0}(1+\Delta)$, $T_{-}=T_{0}(1-\Delta)$ with $T_{0}=0.5$,
$\Delta=0.5$; for (b) $\lambda=0.1$. The inset in (b): the effective coupling
from the nonlinearity $k_{12\rm{eff}}-k_{12}$ as a function of interfacial
linear coupling $k_{12}$.
## 6 Discussion
We use a simplified one-dimensional atomic model to investigate the underlying
physics in the interfacial thermal transport; based on which we find that the
nonlinearity can enhance the thermal transport and the nonlinearity at
interface can induce the interfacial thermal rectification. For a general
chemical bonding strength across interface, we could represent it as a linear
coupling part and nonlinear coupling part; therefore our model can be applied.
However, to compare with experimental results on real materials, the
simplified model need to be generalized to two-dimensional or three
dimensional with input of the real parameters.
The QSCMF theory is a kind of mean-field approximation based on the
nonequilibrium Green’s function, which is similar as the mean field approach
in the literatures [28, 29]. The QSCMF is a good candidate to solve the
nonlinear problem with arbitrary nonlinearity and arbitrary system-bath
coupling. By comparing to the quantum master equation which is limited by
weak-system-bath coupling, we find the QSCMF is a quite accurate method while
it can be applied to strong system-bath coupling. For the interface problem
with two-layer atoms, the QSCMF has a very high accuracy in the wide range of
temperature and nonlinearity. However, if more layer-atoms includes, the
accuracy decreases. If the QSCMF is generalized to two or three dimensional
interface with more atoms in the interface, the self-consistent process is
time consuming, and we also need to pay more attention to its convergency. The
QSCMF is quite good for two-layer-atoms interface which is a reasonable
approximation for short-range inter-atom coupling; however for long-range
interaction, more layers need to be considered at the interface and the
numerical calculation on the nonlinear thermal transport is challenging.
## 7 Conclusion
Based on the NEGF approach, an efficient QSCMF theory is developed to study
the nonlinear interfacial thermal transport. We find that the nonlinearity can
enhance the interfacial thermal transport at weak linear interfacial coupling
while the enhancement vanishes in the strong linear coupling regime. The
enhancement can exist at large nonlinearity where the effective coupling is
less than the harmonic average of the spring constants of the two semi-
infinite chains. Although the leads are linear, the interfacial nonlinearity
can induce rectification provided that the two leads are asymmetric. With
increasing the nonlinear coupling, the rectification first increases then
decreases. The interfacial rectification also depends on the interfacial
linear coupling; it vanishes if the linear interfacial coupling increases far
beyond the nonlinear coupling.
## 8 Acknowledgements
L. Z. and B. L. are supported by the grant R-144-000-300-112 from Ministry of
Education of Republic of Singapore. J.-S. W. and J. T. acknowledge support
from a faculty research grant R-144-000-257-112 of NUS. D. H. surported by
NSFC (Grant No. 11047185) and FRFCU (Grant No. 2010121009) of China.
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|
arxiv-papers
| 2013-09-12T14:10:26 |
2024-09-04T02:49:50.824192
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Lifa Zhang, Juzar Thingna, Dahai He, Jian-Sheng Wang, and Baowen Li",
"submitter": "Lifa Zhang",
"url": "https://arxiv.org/abs/1309.3163"
}
|
1309.3166
|
# Thermal transport across metal-insulator interface via electron-phonon
interaction
Lifa Zhang,1 Jing-Tao Lü,2 Jian-Sheng Wang,1 Baowen Li1,3,4 1 Department of
Physics and Centre for Computational Science and Engineering, National
University of Singapore, Singapore 117542, Republic of Singapore 2 DTU-
Nanotech, Department of Micro- and Nanotechnology, Technical University of
Denmark (DTU), Ørsteds Plads, Bldg. 345E, DK-2800 Lyngby, Denmark 3 NUS
Graduate School for Integrative Sciences and Engineering, Singapore 117456,
Republic of Singapore 4 NUS-Tongji Center for Phononics and Thermal Energy
Science and Department of Physics, Tongji University, 200092 Shanghai, PR
China
###### Abstract
The thermal transport across metal-insulator interface can be characterized by
electron-phonon interaction through which an electron lead is coupled to a
phonon lead if phonon-phonon coupling at the interface is very weak. We
investigate the thermal conductance and rectification flowing between the
electron part and the phonon part using nonequilibrium Green’s function
method. It is found that the thermal conductance has a nonmonotonic behavior
as a function of average temperature or the coupling strength between the
phonon leads in the metal part and the insulator one. The metal-insulator
interface shows evident thermal rectification effect, which can reverse with
changing of average temperature or the electron-phonon coupling.
###### pacs:
68.35.-p, 66.70.-f, 63.20.kd, 72.10.Di
## 1 Introduction
With the increase of integration density, accumulation of heat becomes a
bottleneck for further development of microelectronic devices. Moreover, as
most of electronic devices consist of metal and insulator/semiconductor
interfaces, to understand the thermal transport through the metal-
insulator/semiconductor interface is indispensable for heat dissipation [1, 2,
3].
Figure 1: (a) Heat generation in a metal with an applied bias. The arrow shows
the heat transport from the metal to the insulating substrate. (b) An electron
bath and a phonon bath at two different temperatures connected by EPI. Energy
exchanges between the two baths and there is a net thermal current $J$ from
the hot bath to the cold one. (c) The lattice model of the metal-insulator
interface. The one dimensional electron lead (the left upper semi-infinite
chain where the dots represent electron states ) connects to the one
dimensional atomic chain (the right semi-infinite chain) via electron-phonon
interaction $V_{\rm ep}$. The phonon part in the metal (the left lower semi-
infinite atomic chain) can also be considered to the thermal transport, which
connects to the right atomic chain by a relative coupling $\lambda$.
In metals electrons dominate thermal transport while phonons do in
semiconductors and insulators; thus for thermal transport across metal-
nonmetal interfaces, energy must transfer between electrons and phonons [4].
Thanks to many remarkable physical properties and prospects for large-area
epitaxial growth, graphene is a promising material for future electronics [5].
The biased graphene on a dielectric insulator often appears in electronic
devices especially in transistors. Due to the bias the electrons in metallic
part have a higher effective temperature than the phonons in the insulating
substrate; and the thermal transport from the electrons in metal to the
phonons in insulator is one important channel of energy transfer [7, 8]. It is
found in a very recent experiment on a carbon nanotube that more that 80
percent of the electrical power supplied to the nanotube is transferred
directly into the substrate by the electron-phonon interaction (EPI) [9].
Therefore to understand the thermal transport across the interface through the
EPI is highly desirable for the heat dissipation in electronic devices.
To study the thermal transport across interface, the acoustic mismatch model
[6] and the diffuse mismatch model [10, 4] are widely applied; however, both
models offer limited accuracy in nanoscale interfacial resistance predictions
[11] because they neglect atomic details of actual interfaces. A scattering
boundary method within the lattice dynamic approach [12, 13, 14] fully
considers the atomic structures in the interface; but it can only be applied
to ballistic thermal transport. Classical molecular dynamics is another widely
used method in phonon transport [15, 16, 17], which is not accurate below the
Debye temperature, and ignores the quantum effect. Only recently the
nonequilibrium Green’s function method, which originates from the study of
electronic transport [18], been applied to study the quantum phonon transport
[19, 20, 21, 22]. So far the study of the coupled electronic and phonon
transport [23, 24, 25] is rare, especially in the metal-nonmetal interface [4,
28, 29].
In this paper, using the nonequilibrium Green’s function method, we study the
thermal transport across metal-insulator interface via the EPI. Our model can
also be applied to metal-semiconductor interface. We study energy flow between
an electron bath and a phonon bath which are connected by electron-phonon
coupling, as shown in Fig.1 (b). The electron bath, described by a semi-
infinite electronic chain under tight-binding approximation, connects with a
phonon lead illustrated as a semi-infinite harmonic atomic chain by a weak
electron-phonon interaction, as shown in Fig.1 (c), where the phonon degrees
of freedom in metal can also be considered in the thermal transport.
## 2 Model and Method
We study the interfacial thermal transport at the metal-insulator interface as
shown in Fig. 1(a). To manifest the effect of the interface we exclude the
nonlinear electron or phonon transport in the two materials themselves, thus
the only thermal resistance comes from the interface. Such model really
uncovers the thermal transport properties of the interface itself. To study
the longitudinal thermal transport, that is, the cross-plane interfacial
transport, we can simplify the problem further to thermal transport between an
electron bath and a phonon bath which are connected by an electron-phonon
coupling, as shown in Fig. 1(b), and can be represented by one dimensional
lattice model as shown in Fig. 1(c), where a semi-infinite electronic chain
(at temperature $T_{Le}$) connects to a semi-infinite atomic chain (at
temperature $T_{R}$) by an electron-phonon interaction $V_{ep}$. The phonon
part in the metal can also be considered by another semi-infinite atomic chain
which connects to the right atomic chain by a relative coupling $\lambda$. In
our model, the scattering for both electrons and phonons only comes from the
interface, and the electron and phonon transport in the corresponding semi-
infinite periodic leads is ballistic. Therefore we can partition the system
into three parts ($L$, $C$, $R$), the atoms at the interface are regarded as
center illustrated as the dashed-line rectangle in Fig.1 (c), and the rest
parts are leads. For simplicity, we only consider one dimensional case,
ignoring electron spin. Thus the Hamiltonian of the whole system in Fig.1 (c)
is written as
$H=H^{L}_{\rm e}+H^{C}+H^{LC}_{\rm e}+\sum_{\alpha=L,R}(H^{\alpha}_{\rm
p}+H^{\alpha C}_{\rm p}),$ (1)
where $H^{L}_{\rm
e}=\sum_{i}\varepsilon_{0}c^{\dagger}_{i}c_{i}-\sum_{|i-j|=1}tc^{\dagger}_{i}c_{j}$
is the Hamiltonian of the electron lead in the left part. The electron
coupling between the left lead (the site 1) and the center is $H^{LC}_{\rm
e}=-tc^{\dagger}_{1}c-tc^{\dagger}c_{1}$. We set atomic mass $m=1$, and the
Hamiltonian of the center is
$H^{C}=\frac{1}{2}\dot{u}^{2}+\frac{1}{2}(\lambda
k^{L}+k^{R})u^{2}+\varepsilon_{0}c^{\dagger}c+V_{\rm ep}c^{\dagger}cu,$ (2)
where $k^{L}$ and $k^{R}$ are the spring constants of the left and right
atomic chains, respectively. $u$ is the atom displacement. $\lambda k^{L}$
denotes the coupling between the two atomic chains; $\lambda$ can be chosen
from 0 to 1. $V_{\rm ep}$ is the electron-phonon coupling. The phonon
Hamiltonian of two leads and its coupling to the center are $H^{\alpha}_{\rm
p}=\frac{1}{2}\sum_{i}\dot{u}^{\alpha}_{i}\dot{u}^{\alpha}_{i}+\frac{1}{2}\sum_{|i-j|=0,1}u^{\alpha}_{i}K^{\alpha}_{ij}u^{\alpha}_{j}$
and $H^{\alpha C}_{\rm p}=u^{\alpha}_{1}K^{\alpha C}u$ [20, 23].
Applying the standard procedure of nonequilibrium Green’s function method
[20], without the EPI, we obtain the electron retarded Green’s function
$G^{r}_{0}(\varepsilon)=\left[\varepsilon-\varepsilon_{0}-\Sigma^{r}_{L}(\varepsilon)\right]^{-1}$,
where $\Sigma^{r}_{L}=g^{r}_{L}t^{2}$ is the retarded self-energy with a
surface Green’s function $g^{r}_{L}$, where
$g^{r}=\left[(\varepsilon+i\eta)I-H^{L}_{\rm e}\right]^{-1}$. The less Green’s
function can also be easily obtained. The phonon retarded Green’s functions is
$D^{r}_{0}(\omega)={D^{a}_{0}}^{\dagger}(\omega)=\left[\omega^{2}-(\lambda
k^{L}+k^{R})-\Pi^{r}_{L}(\omega)-\Pi^{r}_{R}(\omega)\right]^{-1}$, where
$\Pi^{r}_{L}(\omega)=(\lambda k^{L})^{2}d^{r}_{L}(\omega)$ and
$\Pi^{r}_{R}(\omega)=(k^{R})^{2}d^{r}_{R}(\omega)$ are the retarded self
energies of left and right leads with surface Green’s function $d^{r}_{L}$ and
$d^{r}_{R}$, where
$d^{r}_{L,R}=\left[(\omega+i\eta)^{2}I-K^{L,R}\right]^{-1}$. The EPI is
included as perturbation. The full Green’s functions are obtained from the
Dyson equation, that is, $G^{r,a}=[(G_{0}^{r,a})^{-1}-\Sigma^{r,a}_{\rm
ep}]^{-1}$, and $G^{<,>}=G^{r}(\Sigma^{<,>}_{\rm ep}+\Sigma^{<,>}_{L})G^{a}$,
where $\Sigma^{r,a,<,>}_{\rm ep}$ is the self-energy from the EPI. For phonon,
we have $D^{r,a}=[(D_{0}^{r,a})^{-1}-\Pi^{r,a}_{\rm ep}]^{-1}$, and
$D^{<,>}=D^{r}(\Pi^{<,>}_{\rm ep}+\Pi^{<,>}_{L}+\Pi^{<,>}_{R})D^{a}$, where
the nonlinear self energy $\Pi^{r,a,<,>}_{\rm ep}$ comes from the electron-
phonon coupling. Keeping the lowest non-zero order (second order) of the self-
energies, we could obtain the nonlinear self-energies $\Sigma^{r,a,<,>}_{\rm
ep}$ and $\Pi^{r,a,<,>}_{\rm ep}$ up the second order, which are general for
any dimensional systems. For our one-dimensional simple model as shown in
Eq.2, the self-energies could be written as
$\Sigma_{\rm ep}^{>,<}(\varepsilon)=iV_{\rm
ep}^{2}\int{G_{0}^{>,<}}(\varepsilon-\omega){D_{0}^{>,<}}(\omega)\frac{d\omega}{2\pi},$
(3)
and
$\displaystyle\Sigma_{\rm ep}^{r}(\varepsilon)$ $\displaystyle=$
$\displaystyle iV_{\rm
ep}^{2}\bigl{\\{}-{D_{0}^{r}}(\omega^{\prime}=0)\int{G_{0}^{<}}(\varepsilon^{\prime})\frac{d\varepsilon^{\prime}}{2\pi}$
(4)
$\displaystyle+\int\frac{d\omega}{2\pi}\bigr{[}{G_{0}^{r}}(\varepsilon-\omega){D_{0}^{<}}(\omega)+{G_{0}^{<}}(\varepsilon-\omega){D_{0}^{r}}(\omega)$
$\displaystyle+{G_{0}^{r}}(\varepsilon-\omega){D_{0}^{r}}(\omega)\bigr{]}\bigr{\\}}.$
The nonlinear self-energies for the phonons are
$\Pi_{\rm ep}^{>,<}(\omega)=-iV_{\rm
ep}^{2}\int\frac{d\varepsilon}{2\pi}{G_{0}^{>,<}}(\varepsilon){G_{0}^{<,>}}(\varepsilon-\omega),$
(5)
and
$\displaystyle\Pi_{\rm ep}^{r}(\omega)$ $\displaystyle=$ $\displaystyle-
iV_{\rm
ep}^{2}\int\frac{d\varepsilon}{2\pi}\bigl{[}{G_{0}^{r}}(\varepsilon){G_{0}^{<}}(\varepsilon-\omega)$
(6)
$\displaystyle+{G_{0}^{<}}(\varepsilon){G_{0}^{a}}(\varepsilon-\omega)\bigr{]}.$
Eqs. (3-6) are the so-called Born approximation (BA). By replacing the bare
Green’s functions $G_{0}$ and $D_{0}$ with the full Green’s functions $G$ and
$D$, we do iteration under self-consistent Born approximation (SCBA). While
the BA fails to satisfy the energy current conservation, the SCBA fulfills it
[23].
Therefore, we obtain the heat current from electron part as [18, 30, 31]
$J^{\rm e}=\int\frac{d\varepsilon}{2\pi\hbar}\varepsilon~{}{\rm
Tr}\\{G^{>}(\varepsilon)\Sigma^{<}_{L}(\varepsilon)-G^{<}(\varepsilon)\Sigma^{>}_{L}(\varepsilon)\\}.$
(7)
The heat current from the phonon lead is [20, 23, 32]
$J^{\rm p}_{\alpha}=-\int\frac{d\omega}{4\pi}\hbar\omega~{}{\rm
Tr}\\{D^{>}(\omega)\Pi^{<}_{\alpha}(\omega)-D^{<}(\omega)\Pi^{>}_{\alpha}(\omega)\\}.$
(8)
Due to energy conservation, $J^{\rm e}+J^{\rm p}_{L}+J^{\rm p}_{R}=0$. If the
coupling $\lambda=0$, we obtain $J^{\rm e}=-J^{\rm p}_{R}$ when the energy
flows across the interface only through EPI, that is, the heat generation of
the metal part transfers directly to the insulator part. We set the Plank
constant $\hbar=1$ and the Boltzman constant $k_{\rm B}=1$ in the following
numerical calculation; and $\varepsilon_{0}=0$, $k^{L}=k^{R}=1,\;t=1$. We
define the conductance $\sigma=|J^{\rm e}|/\Delta T$ to illustrate the thermal
transport between the electron part and phonon part.
## 3 Thermal Conductance across the Interface
Figure 2: (a) Thermal conductance across the metal-insulator interface vs the
average temperature of the baths. $T_{L\rm e}=T_{0}+\Delta T$, $T_{L\rm
p}=T_{R}=T_{0}-\Delta T$, $\Delta T=0.01$, $V_{\rm ep}=0.01$. (b) The thermal
current density $J(\varepsilon)$ ($J^{\rm e}=\int d\varepsilon
J(\varepsilon)$) for electrons with energy $\varepsilon$ at
$\lambda=0^{+}=10^{-4}$. (c) The thermal conductance (left scale) and current
(right scale) from the electron bath vs the temperature difference $\Delta T$
at $T_{0}=1.0$.
As shown in Fig.1(a), if the metal part is applied by a voltage, the average
energy of electrons increases, which means the electrons have a higher
effective temperature than the phonons. Thus we set the left and right phonon
leads to a lower temperature $T_{L\rm p}=T_{R}=T_{0}-\Delta T$, while the
electron lead has a higher temperature $T_{L\rm e}=T_{0}+\Delta T$. If we
don’t consider the phonon contribution in the metal part, the thermal
conductance from the electron lead to the phonon lead has a nonmonotonic
behavior with increasing the average temperature $T_{0}$, as shown in Fig.
2(a). In the low-temperature region, as temperature increases, the thermal
conductance increases due to more contribution from electrons with energy far
from Fermi surface ($\varepsilon=0$) although there is certain decrease from
the contribution near Fermi surface, which can be seen from the solid line
($T_{0}=0.2$) and dashed line ($T=0.6$) in Fig. 2 (b). With the temperature
further increasing, the dominating effect comes from the decrease of the
contribution near the Fermi level, as shown in the dashed line and dotted line
of Fig. 2 (b). This causes the thermal conductance to decrease after certain
temperature. Such dependence of thermal conductance on temperature is
consistent with the experiment observation of thermal contact conductance
between graphene and silicon dioxide [33]. While the conductance has a
nonmonotonic dependence on $T_{0}$, the heat current and conductance shows a
monotonic behavior with temperature difference $\Delta T$ increasing, as shown
in the inset of Fig. 2 (c). If we include the phonon contribution in the metal
part, the thermal conductance has a similar temperature dependence, but with a
smaller magnitude. This means that the heat generation from the electron part
is affected by the coupling strength between the two phonon parts in the
metal-insulator interface.
In order to investigate the effect of phonon-phonon channel $\lambda$ on the
thermal transport through electron-phonon channel $V_{ep}$, in Fig. 3 (a) we
plot the heat current from the electrons to the phonons when we increase the
coupling strength $\lambda$ between the two phonon parts. We find that in the
small $\lambda$ region, with $\lambda$ increasing the heat current flows into
the right phonon lead will increase due to the smaller difference between the
atom ($K^{C}=(\lambda k^{L}+k^{R})$) in interface and the atoms
($K^{R}_{ii}=2k^{R}$) in the right lead, thus the scattering for phonons at
the interface decreases. The heat is easier to transport to the right lead;
$J^{\rm e}$ increases. At the same time, the heat can flow from the electrons
to the left phonon lead; $-J^{\rm p}_{L}$ increases with $\lambda$. When the
$\lambda$ increases further, the phonon scattering between left chain and the
center atom decreases, more thermal energy transferred from electron is easier
to flow into the left lead. Such fast increasing of $-J^{\rm p}_{L}$ causes
the decreasing of the heat into the right lead $-J^{\rm p}_{R}$. With
increasing of $\lambda$, the phonon scattering at the interface between the
left and right atomic chains decreases, more phonons can transport coherently
in the two atomic leads, and are less efficient to couple to electrons; thus
the heat flowing from the electron lead begins to decrease. Due to continuous
increase of the heat flowing into the left phonon lead $-J^{\rm p}_{L}$, and
$J^{\rm e}+J^{\rm p}_{L}+J^{\rm p}_{R}=0$, thus the decrease of heat from
electrons $J^{\rm e}$ lags behind the decrease of the heat into the right lead
$-J^{\rm p}_{R}$. And if $\lambda=1$, the heat into two phonon leads has the
same value due to the symmetry. As shown in Fig. 3 (b), if the non-biased
metal part in a higher temperature while the insulator in a lower one, the
heat flow from electron part has a similar dependence on $\lambda$; but the
heat flowing from left phonon lead and the one into the right phonon lead
monotonically increase with $\lambda$ since phonon transport dominates heat
flow in such case and the phonon scattering decreases. We also find that, in
this case the thermal transport through the phonon-phonon channel $\lambda$ is
several orders of magnitude larger than that through the electron-phonon
channel. And even if $\lambda=4\times 10^{-4}$ (see the inset of Fig. 3(b)),
the thermal current of phonon is ten times larger than that of electron.
Therefore for the thermal transport in the non-biased metal-insulator
interface, such small thermal current of electron can be ignored, which is
consistent with the recent finding in Ref.[34]. However, for the interface
between a biased metal and a insulator, the thermal transport through the
electron-phonon interaction is very important, which dominated the thermal
transport even if $\lambda=1$ as shown in Fig 3 (a), which can explain the
remote Joule heating via electron-phonon interaction in Ref. [9].
Figure 3: (a) Thermal current vs the coupling $\lambda$. The solid line
corresponds to the thermal current $J^{\rm e}$ flowing from the electron bath.
The dashed and dotted lines correspond to the thermal current ($-J^{\rm
p}_{L}$ and $-J^{\rm p}_{R}$ ) flowing into the left and right phonon baths,
respectively. (b) Thermal current vs the coupling $\lambda$ at $T_{L\rm
e}=T_{L\rm p}=1.01$, $T_{R}=0.99$. The solid, dashed and dotted lines
correspond to the thermal current flowing from the electron bath (left scale),
the left and the right phonon baths (right scale), respectively. The inset is
the zoom-in of (b) for small $\lambda$. For (a) and (b), $V_{\rm ep}=0.01$.
## 4 Thermal Rectification across the Interface via EPI
Figure 4: Thermal rectification of the metal-insulator interface vs
temperature difference. The inset: the thermal current across the interface vs
temperature difference. $T_{h}=T_{0}(1+\Delta T)$, $T_{c}=T_{0}(1-\Delta T)$,
$V_{\rm ep}=0.01$.
Thanks to nonlinearity and asymmetry of the metal-insulator interface, we can
expect the thermal rectification, which is defined as $R=(J_{+}-J_{-})/{\rm
max}\\{J_{+},J_{-}\\}$, where $J_{+}$ is the forward direction heat flux,
defined as $T_{L}=T_{h},T_{R}=T_{c}$, and $J_{-}$ is that of the backward
direction when $T_{L}=T_{c},T_{R}=T_{h}$. Here, $T_{h}$ and $T_{c}$ correspond
to the temperatures of the hot and cold baths, respectively. In the rest of
the paper, we will not consider the phonon contribution in metal (we set
$\lambda=10^{-4}$, $T_{Lp}=T_{R}$ in the calculation to avoid divergence),
which will not change the physical properties of the rectification, and only
decrease its magnitude. Figure 4 shows the dependence of thermal rectification
on temperature difference for different average temperatures. The
rectification shows a monotonically increasing behavior of temperature
difference due to the bigger difference of heat current in larger temperature
difference as shown in the inset of Fig. 4, which is consistent with all the
traditional study on thermal rectification [35, 36]. With the temperature
$T_{0}$ increasing, we find that the rectification can change sign.
Figure 5: (a) Thermal rectification of the metal-insulator interface vs
temperature $T_{0}$ for different temperature differences at $V_{\rm
ep}=0.01$. (b) Thermal rectification vs the electron-phonon interaction
$V_{ep}$ at different temperatures. (c) The forward thermal current (solid
line) and backward thermal one (dashed line) vs $V_{\rm ep}$ at $T_{0}=0.25$.
(d1)-(d4) The local available density of states (LADOS) at surface for the
left electron part ($\rho_{\rm e}={|\rm Img}(g^{<})|$)(solid line) and the
right phonon part ($\rho_{\rm p}={|\rm Img}(d^{<})\omega|$)(dashed line).
(d1)((d2)) and (d3)((d4)) are corresponding to the forward (backward)
transport at $T_{0}=0.25$ and $1.0$, respectively. (e1)-(e4) The forward
(solid lines) and backward (dashed lines) thermal current density
$J(\varepsilon)$ for different temperatures and different EPI. For all the
curves, $T_{h}=T_{0}(1+\Delta T)$, $T_{c}=T_{0}(1-\Delta T)$; from (b) to (e)
$\Delta T=0.4$.
The dependence of thermal rectification on temperature and the EPI is shown in
Fig. 5. We find that at a relative larger EPI ($V_{\rm ep}=0.01$), with
increasing the temperature, the thermal rectification can change sign from
positive to negative. However for a very weak EPI, the rectification always
keeps negative, as shown in Fig 5(b). Within the range of $\lambda=10^{-4}\sim
10^{-2}$, the rectification is always negative for higher temperatures, while
it can change sign from negative to positive in a lower temperature $T=0.25$.
The relation of thermal currents in the forward and backward direction changes
with increasing of $\lambda$ while both their magnitudes monotonically
increase, as shown in Fig 5(c).
From the local available density of states (LADOS) of electrons and phonons at
the surface of the corresponding leads, we find that at a higher temperature
$T_{0}=1.0$, going from the forward transport to the backward one, the main
change is that the LADOS phonons largely increases, as shown in Fig. 5 (d3)
and (d4), which causes a larger current in the backward direction than that in
the forward direction, thus the rectification is negative. At a lower
temperature, the change of LADOS for electrons and phonons are not obvious, as
shown in Fig. 5 (d1) and (d2), and the rectification could change with the
strength of EPI.
Turning on the EPI, the electron Green’s functions change much more than the
phonon ones; with increasing EPI, the Green’s functions of electrons $G^{>}$
and $G^{<}$ become more asymmetric, especially for the $G^{<}$ which have
larger value at the energy far away from the Fermi level ($\varepsilon=0$). At
a very weak EPI, the heat current comes from the contribution of electronic
energies symmetrically away from the Fermi level. The heat current of the
backward transport is larger than that of the forward one, as shown in Fig. 5
(e1) and (e3), due to the increase of the LADOS of phonons in the backward
direction, thus the rectification are negative for all the temperatures at
weak EPI as shown in Fig 5(b). With a stronger EPI, which induces larger value
of $G^{<}$ at the energy range $-2\sim-1$, where the center electron has
larger density of states. Thus we find that the forward current has larger
value at the range of $-2\sim-1$ as shown in Fig. 5(e2). At the same time, the
heat current can be affected by the temperature, in the low temperature range,
as temperature increases, more electrons with energy far away from Fermi level
contribute to the thermal current, as discussed in Fig. 2(b). Thus in Fig.
5(e2), the forward heat current is larger than the backward one; the
rectification is positive at low temperatures with larger EPI. In the higher
temperature range, the contribution from energy far away from Fermi level will
decrease with temperature increasing as shown in Fig. 2(b); further more, the
mainly change from LADOS of electron and phonon is the larger phonon LADOS in
the backward transport than that in the forward one; thus the backward heat
current is larger than the forward one as shown in Fig. 5(e4). Therefore, due
to mismatch of electron LADOS and phonon one, the rectification is negative
for high temperatures or weak EPI. Due to the nonlinearity of EPI and its
relation with temperature, we can find the positive rectification with
relative larger EPI at low temperatures.
## 5 Discussions
Figure 6: Green’s functions for electrons and phonons, $T_{L\rm
e}=T_{0}+\Delta T$, $T_{L\rm p}=T_{R}=T_{0}-\Delta T$, $T_{0}=0.2$, $\Delta
T=0.01$. (a) The imagine part of the bare greater and less Green’s functions
for phonons and electrons: $D_{0}^{>}$ (solid line),$D_{0}^{<}$ (dotted
line),$G_{0}^{>}$ (dashed line),$G_{0}^{<}$ (dash-dotted line). (b) The
imagine part of the full Green’s functions for phonons and electrons: $D^{>}$
(solid line),$D^{<}$ (dotted line),$G^{>}$ (dashed line),$G^{<}$ (dash-dotted
line) with the EPI ($V_{\rm ep}=0.01$).
### 5.1 Full Green’s functions in the calculation
Under the SCBA, we calculate the self energies of the EPI; after the iteration
is convergent we obtain all the Green’s functions of $G^{>,<}$ and $D^{>,<}$
in Eq. (3) and (4) in the main text, then we can calculate the heat currents.
If we do not consider the phonon contribution in the metal part, we should
choose the coupling $\lambda$ be a very tiny positive value, thus the heat
current to the left part will be very small and negligible; as shown in Fig. 3
in the main text, if $\lambda=10^{-4}$ the heat current flowing into the left
phonon leads is almost zero. If $\lambda=0$, the iteration to calculate self
energies for EPI will be divergent, but if $\lambda=0^{+}=10^{-4}$, it can be
convergent. Thus here the tiny $\lambda=0^{+}$ plays a role as very small
onsite potential for the center atom, which is similar as the tiny onsite in
computing the phonon Hall effect [38]. We plot the bare Green’s functions
$G_{0}$, $D_{0}$ and the full Green’s functions $G$, $D$ in Fig. 6. Without
the EPI, the bare Green’s functions of electrons $G_{0}^{<}$ and $G_{0}^{>}$
are symmetric with respect to the origin as shown in Fig. 6 (a), but which
symmetry for $G^{<}$ and $G^{>}$ will be broken and change much if we turn on
the EPI, which can be seen in Fig 6 (b). However, the Green’s functions for
phonons do not change much.
### 5.2 The role of Fermi energy of electrons
Figure 7: The conductance (left panel) and the rectification (right panel) vs
the onsite $\varepsilon_{0}$. $\lambda=10^{-4}$. For comparison, The dotted
lines in (a) and (b) are copied from the corresponding line in Fig.2 (a) and
Fig.5(a), respectively.
We set the onsite $\varepsilon_{0}=0$ in the main text, thus the Fermi energy
of the electron system is $\varepsilon=0$. However we can change the onsite
$\varepsilon$, as shown in Fig. 7, the onsite potential can change the
magnitude of the thermal conductance and the rectification, and the
temperature where the thermal conductance and rectification have maximum value
shifts a bit. However the curves of the thermal conductance and rectification
have the same properties with the case at $\varepsilon_{0}=0$. Therefore, the
mechanism found for the non-monotonicity of conductance and reversal of
rectification still holds for the case of nonzero Fermi energy.
### 5.3 Limitation and application
We use a simplified one-dimensional model to study thermal conductance and
rectification across the metal-insulator interface via electron-phonon
interaction. For some quasi-one dimensional metal-insulator interface systems,
the thermal transport can simply be recast to one-dimensional model, then our
results of the conductance and rectification can be applied. The thermal
resistance of a whole interface system comes from the two materials themselves
and the interface between them. Our study in this work focuses on the
interface itself, which provides understanding on the interface in the thermal
transport. Although the thermal transport via electron-phonon interaction can
be ignored in the non-biased metal-insulator interface, it dominates the
thermal transport in the remote Joule heating in a biased metal-insulator
interface.
For real materials, there would be more atoms in one unit cell, thus more
phonon branches will contribute to the thermal transport. To compare with
experiments on real materials, the simplified model need to be generalized to
two-dimensional or three dimensional with input of the real parameters. For
realistic three dimensional systems, the density of states of electrons and
phonons are quite different from those in one dimension [39], then the
analysis for the mechanism of the thermal rectification would be different and
more complicated; but the general properties on the interfacial thermal
transport, such as thermal conductance as functions of $T$, $\lambda$ and the
rectification tuned by temperature or electron-phonon interaction, would still
hold and could be verified by the future experiments.
The nonequilibrium Green’s function approach is a good candidate to solve the
interfacial thermal transport. If we study two or three dimensional systems
with large number atoms in the interface, the self-consistent process is time
consuming, and we also need to pay more attention to its convergency. If we
only are interested on the thermal conductance accurate to the second order of
$V_{ep}$, the Born approximation is applicable to avoid the convergent
problem. However, for the thermal rectification as a function of $V_{ep}$ can
only be obtained by the self consistent Born approximation. Therefore the
calculation of the simple model can provide the understanding of the basic
properties of thermal transport at metal-insulator interface and can be
generalized to real materials with the parameters obtained from first-
principle calculation. The challenges would come from the time consuming for
the self-consistent calculation of the nonlinear self energies $\Sigma_{\rm
ep}$ and $\Pi_{\rm ep}$.
## 6 Conclusion
In summary, using nonequilibrium Green’s function method we study the thermal
transport across metal-insulator interface with an EPI. The thermal
conductance and rectification across the interface is thoroughly investigated
under SCBA by a clean and efficient model where an electron lead couples to a
phonon lead through the EPI. We find that the thermal conductance is a
nonmonotonic function of average temperature, it has a maximum value at
certain temperature. Considering the phonon part in metal contributing to the
thermal transport, the heat flow from electron part will change while it has a
similar dependence on temperature. The thermal rectification effect across the
metal-insulator interface can reverse with varying of system average
temperature and the EPI. Our results are very helpful to explain and guide the
experiment on the thermal transport and dissipation in electronic devices and
interface system; and could be verified in the biased and non-biased graphene
sheet supported on insulating substrate. The reversal rectification is
significant for the study of one-way heat transport and tuning the direction
of thermal transport; thus it could have wide application in the energy
science.
## Acknowledgements
L. Z.and B. L. are supported by the grant R-144-000-300-112 from Ministry of
Education of Republic of Singapore. J.-S. W. acknowledges support from a URC
research grant R-144-000-257-112 of NUS. J.-T. L. acknowledges the Lundbeck
Foundation for financial support (R49-A5454).
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|
arxiv-papers
| 2013-09-12T14:20:03 |
2024-09-04T02:49:50.831822
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Lifa Zhang, Jing-Tao L\\\"u, Jian-Sheng Wang, Baowen Li",
"submitter": "Lifa Zhang",
"url": "https://arxiv.org/abs/1309.3166"
}
|
1309.3277
|
11affiliationtext: Department of Astronomy, University of Texas at Austin,
Austin, TX 78712, USA22affiliationtext: Racah Institute for Physics, The
Hebrew University, Jerusalem, 91904, Israel
# Magnetic Fields In Relativistic Collisionless Shocks
Rodolfo Santana1a1aaffiliationmark: , Rodolfo Barniol
Duran2b2baffiliationmark: , & Pawan Kumar1c1caffiliationmark: (a)
[email protected]; (b) [email protected] (c)
[email protected]
###### Abstract
We present a systematic study on magnetic fields in Gamma-Ray Burst (GRB)
external forward shocks (FSs). There are 60 (35) GRBs in our X-ray (optical)
sample, mostly from Swift. We use two methods to study $\epsilon_{B}$
(fraction of energy in magnetic field in the FS). 1. For the X-ray sample, we
use the constraint that the observed flux at the end of the steep decline is
$\geq$ X-ray FS flux. 2. For the optical sample, we use the condition that the
observed flux arises from the FS (optical sample light curves decline as $\sim
t^{-1}$, as expected for the FS). Making a reasonable assumption on $E$ (jet
isotropic equivalent kinetic energy), we converted these conditions into an
upper limit (measurement) on $\epsilon_{B}n^{2/(p+1)}$ for our X-ray (optical)
sample, where $n$ is the circumburst density and $p$ is the electron index.
Taking $n=1\mbox{ cm}^{-3}$, the distribution of $\epsilon_{B}$ measurements
(upper limits) for our optical (X-ray) sample has a range of $\sim
10^{-8}-10^{-3}$ ($\sim 10^{-6}-10^{-3}$) and median of $\sim\mbox{ few}\times
10^{-5}$ ($\sim\mbox{ few}\times 10^{-5}$). To characterize how much
amplification is needed, beyond shock compression of a seed magnetic field
$\sim 10\mu\mbox{G}$, we expressed our results in terms of an amplification
factor, $AF$, which is very weakly dependent on $n$ ($AF\propto n^{0.21}$).
The range of $AF$ measurements (upper limits) for our optical (X-ray) sample
is $\sim 1-1000$ ($\sim 10-300$) with a median of $\sim 50$ ($\sim 50$). These
results suggest that some amplification, in addition to shock compression, is
needed to explain the afterglow observations.
###### Subject headings:
gamma-ray burst: general -radiation mechanisms: non-thermal - methods:
analytical
## 1\. Introduction
Gamma-Ray Bursts (GRBs) are bright explosions occurring at cosmological
distances which release gamma-rays for a brief time, typically on a timescale
of $\sim\mbox{ few}\times$ 10 sec (e.g. Piran 2004, Gehrels et al. 2009, Zhang
2011). This short-lived emission of gamma-rays is known as the prompt
emission. After the prompt emission, long-lived emission in the X-ray,
optical, and radio bands (on timescales of days, months, or even years) is
also observed from what is called the “afterglow”. Although the mechanism for
the prompt emission is currently being debated, the afterglow emission has a
well-established model based on external shocks (Rees & Meszaros, 1992;
Meszaros & Rees, 1993; Paczynski & Rhoads, 1993). In this framework, a
relativistic jet emitted by the central engine interacts with the medium
surrounding the GRB progenitor. This interaction produces two shocks; the
external-reverse shock and the external-forward shock (Meszaros & Rees, 1997;
Sari & Piran, 1999b). The external-reverse shock heats up the jet while the
external-forward shock heats up the medium surrounding the explosion. The
external-reverse shock is believed to be short-lived in the optical band (Sari
& Piran, 1999a) and might have been observed, perhaps, in a few cases. The
long-lived afterglow emission is interpreted as synchrotron radiation from the
external-forward shock. This shock is taken to produce a power-law
distribution of high energy electrons and to amplify the pre-existing seed
magnetic field in the surrounding medium. These high energy electrons are then
accelerated by the amplified magnetic field and emit radiation by the
synchrotron process.
One of the open questions in the field of GRB afterglows is: what is the
dynamo mechanism amplifying magnetic fields in the collisionless relativistic
shocks involved for GRB external shocks? The magnetic field strength
downstream of the shock front is expressed in terms of $\epsilon_{B}$, which
is defined as the fraction of energy that is in the magnetic field downstream
of the shock front. With this definition, the explicit expression for
$\epsilon_{B}$ is
$\epsilon_{B}=\frac{B^{2}}{32\pi nm_{p}c^{2}\Gamma^{2}},$ (1)
where $B$ is the co-moving magnetic field downstream of the shock front, $n$
is the density surrounding the GRB progenitor, $m_{p}$ is the proton mass, $c$
is the speed of light, and $\Gamma$ is the Lorentz factor of the shocked fluid
downstream of the shock front (e.g. Sari et al., 1998; Wijers & Galama, 1999;
Panaitescu & Kumar, 2000). If shock compression is the only mechanism
amplifying the magnetic field downstream of the shock front, then $B$ is given
by $B=4\Gamma B_{0}$ (e.g. Achterberg et al., 2001), where $B_{0}$ is the seed
magnetic field in the medium surrounding the GRB progenitor. Using this
expression for $B$, $\epsilon_{B}$ simplifies to $\epsilon_{B}=B_{0}^{2}/2\pi
nm_{p}c^{2}$. Using the value for the ambient magnetic field of the Milky Way
galaxy $B_{0}\sim\mbox{ few }\mu\mbox{G}$ and a density for the surrounding
medium of $n=1\mbox{ cm}^{-3}$, $\epsilon_{B}$ is expected to be $\sim
10^{-9}$.
Several studies have modeled afterglow data to determine what values of the
afterglow parameters best describe the observations (e.g. Wijers & Galama,
1999; Panaitescu & Kumar, 2002; Yost et al., 2003; Panaitescu, 2005). The
results from previous studies show that $\epsilon_{B}$ ranges from
$\epsilon_{B}\sim 10^{-5}-10^{-1}$. These values for $\epsilon_{B}$ are much
larger than the $\epsilon_{B}\sim 10^{-9}$ expected from shock compression
alone and suggest that some additional amplification is needed to explain the
observations. There have been several theoretical and numerical studies that
have considered possible mechanisms, operating in the plasma in the medium
surrounding the GRB, that can generate extra amplification for the magnetic
field. The mechanisms that have been proposed are the two-stream Weibel
instability (Weibel, 1959; Medvedev & Loeb, 1999; Gruzinov & Waxman, 1999;
Silva et al., 2003; Medvedev et al., 2005) and dynamo generated by turbulence
(Milosavljević & Nakar, 2006; Milosavljevic et al., 2007; Sironi & Goodman,
2007; Goodman & MacFadyen, 2008; Couch et al., 2008; Zhang et al., 2009;
Mizuno et al., 2011; Inoue et al., 2011).
Recent results (Kumar & Barniol Duran, 2009, 2010) found surprisingly small
values of $\epsilon_{B}\sim 10^{-7}$ for 3 bright GRBs with Fermi$/$LAT
detections. These values of $\epsilon_{B}$ are $\sim 2$ orders of magnitude
smaller than the smallest previously reported $\epsilon_{B}$ value and they
can be explained with the only amplification coming from shock compression of
a seed magnetic field of a few $10\mu\mbox{G}$ 111The values given above of
$\epsilon_{B}\sim 10^{-7}$ are under the assumption of $n=1\mbox{ cm}^{-3}$.
It is important to note that when reaching the conclusion that shock
compression provides enough amplification to explain the afterglow data, Kumar
& Barniol Duran (2009, 2010) did not assume a value for $n$. Also, the results
of small $\epsilon_{B}\sim 10^{-7}$ values do not depend on whether or not the
LAT emission is produced by the external-forward shock. These small
$\epsilon_{B}$ values were inferred from the late time X-ray and optical data
and from the constraint that the external-forward shock does not produce flux
at 150 keV that exceeds the observed prompt emission flux at 50 seconds..
Although this seed magnetic field is stronger than the one of the Milky Way
galaxy by about a factor $\sim 10$, seed magnetic fields of a few
$10\mu\mbox{G}$ have been measured before. The seed magnetic fields in the
spiral arms of some gas-rich spiral galaxies with high star formation rates
have been measured to be 20-30 $\mu\mbox{G}$ (Beck, 2011). Seed magnetic
fields as high as $0.5-18$ mG were measured in starburst galaxies by measuring
the Zeeman splitting of the OH megamaser emission line at 1667 MHz (Robishaw
et al., 2008).
Given this disagreement between the recent and previous results, the question
regarding the amplification of magnetic fields in GRB external relativistic
collisionless shocks remains unanswered. The first goal of this study is to
provide a systematic determination of $\epsilon_{B}$ for a large sample of
GRBs by using the same method to determine $\epsilon_{B}$ for each burst in
our X-ray or optical sample. This is the first time such a large and
systematic study has been carried out for $\epsilon_{B}$. Knowing the value of
$\epsilon_{B}$ for large samples will help us determine how much amplification
of the magnetic field is needed to explain the afterglow observations. We
mostly limit our samples to GRBs detected by the Swift satellite, with
measured redshift. In this study, we determine an upper limit on
$\epsilon_{B}$ for our X-ray sample and a measurement of $\epsilon_{B}$ for
our optical sample. We use a new method to determine an upper limit on
$\epsilon_{B}$ with X-ray data, which relies on using the steep decline
observed by Swift in many X-ray light curves. We expect that the observed flux
at the end of the steep decline is larger than the predicted flux from the
external-forward shock. Making reasonable assumptions about the other
afterglow parameters, we are able to convert this constraint into an upper
limit on $\epsilon_{B}$. To determine a measurement of $\epsilon_{B}$ for our
optical sample, we restrict our sample to light curves that show a power law
decline with a temporal decay $\sim 1$ at early times, $\sim 10^{2}-10^{3}$
seconds, as expected for external-forward shock emission. We choose this
selection criteria so that the optical emission is most likely dominated by
external-forward shock emission. Making the same reasonable assumptions for
the other afterglow parameters and using the condition that the observed flux
from the optical light curve is equal to the external-forward shock flux, we
are able to convert this condition into a measurement of $\epsilon_{B}$. We
also applied a consistency check for the bursts that are in common to our
X-ray and optical samples to make sure our results for $\epsilon_{B}$ are
correct. The second goal of this study is to determine how much amplification,
in addition to shock compression, is needed to explain the results for the
$\epsilon_{B}$ upper limits$/$measurements. To quantify how much amplification
beyond shock compression is required by the observations, we also express the
results we found for the $\epsilon_{B}$ upper limits (measurements) for our
X-ray (optical) sample in terms of an amplification factor.
This paper is organized as follows. We begin in Section 2 by presenting a
review of the values previous studies have found for the microphysical
afterglow parameters $\epsilon_{e}$ (the fraction of energy in electrons in
the shocked plasma) and $\epsilon_{B}$. In Section 3 (Section 4), we present
the method we use to determine an upper limit (measurement) on $\epsilon_{B}$
and apply it to our X-ray (optical) sample of GRBs. In Section 5, we use the
GRBs that are in common to both samples to perform a consistency check. We
search for a possible correlation between the kinetic energy of the blastwave
and $\epsilon_{B}$ in Section 6. In Section 7, we write our results for
$\epsilon_{B}$ for our X-ray and optical samples in terms of an amplification
factor, which quantifies how much amplification – beyond shock compression –
is required by the observations. Lastly, in Section 8, we discuss our results
and present our conclusions. The convention we use for the specific flux
$f_{\nu}$, the flux per unit frequency $\nu$, is
$f_{\nu}\propto\nu^{-\beta}t^{-\alpha}$. In this convention, $\beta$ is the
spectral decay index and $\alpha$ is the temporal decay index. For a GRB at a
given redshift $z$, when calculating the luminosity distance to the GRB,
$d_{L}$, we used the Cosmological parameters $H_{0}=71\mbox{
km}/\mbox{sec}/\mbox{Mpc}$, $\Omega_{m}=0.27$, and $\Omega_{\Lambda}=0.73$.
## 2\. Literature Review Of Values Of $\epsilon_{e}$ and $\epsilon_{B}$
The flux observed from the external-forward shock depends on 6 parameters.
These parameters are $E$, $n$, $s$, $\epsilon_{e}$, $\epsilon_{B}$, and $p$.
$E$ is the isotropic equivalent kinetic energy of the jet and $n$ is the
number density of the surrounding medium. The density is taken to be
spherically symmetric and to decrease with $r$ as $n(r)\propto r^{-s}$, where
$s$ is a constant determining the density profile of the surrounding medium
and $r$ is the distance from the center of the explosion. Two cases are
usually considered for the density profile: $s=0$ and $s=2$, which
respectively correspond to a constant density medium and a wind medium. The
microphysical parameters are $\epsilon_{e}$ and $\epsilon_{B}$, where
$\epsilon_{e}$ ($\epsilon_{B}$) is the fraction of energy in the electrons
(magnetic field) in the shocked fluid. The microphysical parameters are taken
to be constant throughout the afterglow emission. A power law distribution of
electrons, $dN_{e}/d\gamma_{e}\propto\gamma_{e}^{-p}$ with
$\gamma_{e}\geq\gamma_{i}$, where $\gamma_{e}$ is the Lorentz factor of the
electrons, $\gamma_{i}$ is the minimum Lorentz factor of the electrons, and
$N_{e}$ is the number of electrons, is assumed to be produced when the
external-forward shock interacts with the surrounding medium. The power-law
index of the electron distribution, $p$, is a constant known as the electron
index.
In practice, it is very difficult to determine the values of the 6 afterglow
parameters. The value of $p$ and the density profile of the surrounding medium
(whether we have a $s=0$ or $s=2$ medium) can be determined from observations
of the afterglow spectral decay and temporal decay of the light curve with the
so-called “closure” relations (e.g. Zhang et al., 2006). The remaining 4
afterglow parameters are more difficult to determine. What is needed to
determine these 4 parameters is observations of the afterglow emission in the
4 different spectral regimes of the synchrotron afterglow spectrum (we will
discuss the afterglow synchrotron spectrum in more detail in Section 3.4). In
practice, most GRBs do not have this wealth of observations. In order to
determine these 4 parameters, previous works have either focused only on
determining the afterglow parameters for bursts with high quality data,
spanning all portions of the synchrotron spectrum, or have applied various
simplifying assumptions.
We performed a literature search for papers that determine values for
$\epsilon_{e}$ and $\epsilon_{B}$ to get an idea of what typical values
previous works have found. Different authors applied different techniques for
finding $\epsilon_{e}$ and $\epsilon_{B}$. When displaying the results from
the literature, we did not discriminate against any method and simply plotted
every value we found. However, we did not consider works that made simplifying
assumptions when determining $\epsilon_{e}$ and $\epsilon_{B}$, such as
equipartition of proton and electron energy ($\epsilon_{e}$). The GRBs for
which we found $\epsilon_{e}$ and $\epsilon_{B}$ values are shown in Table
LABEL:epsilon_e_and_epsilon_B_table in Appendix A. Except for GRB 080928, all
the GRBs in our sample have radio, optical, and X-ray observations, allowing
for a determination of all the afterglow parameters. We included GRB 080928 in
our sample because $\epsilon_{e}$ and $\epsilon_{B}$ were able to be uniquely
determined from optical and X-ray observations (Rossi et al., 2011).
The $\epsilon_{e}$ values we found in the literature for 29 GRBs are shown in
the histogram in the left panel of Figure 1. There is a narrow distribution
for $\epsilon_{e}$; it only varies over one order of magnitude, from $\sim
0.02-0.6$, with very few GRBs reported to have $\epsilon_{e}<0.1$. The mean of
this distribution is 0.24 and the median is 0.22. About 62% of the GRBs in
this sample have $\epsilon_{e}\sim 0.1-0.3$. These results for $\epsilon_{e}$
are also supported by recent simulations of relativistic magnetized
collisionless electron-ion shocks presented in Sironi & Spitkovsky 2011, where
they found $\epsilon_{e}\sim 0.2$. The narrow distribution of $\epsilon_{e}$
values from the literature and the results from recent simulations both show
that $\epsilon_{e}$ does not change by much from GRB to GRB.
The $\epsilon_{B}$ values we found in the literature for 30 GRBs are show in
the histogram in the right panel of Figure 1. Comparing the two histograms in
Figure 1, we can immediately see that there is a much wider range in the
distribution of $\epsilon_{B}$, with $\epsilon_{B}$ ranging from $\sim
3.5\times 10^{-5}-0.33$. A noticeable peak, containing about 24% of the
bursts, is seen in the bin with $-1<\mbox{log}_{10}(\epsilon_{B})<-0.5$. Two
other peaks, each containing about 17% of the GRBs, are seen in the bins with
$-2<\mbox{log}_{10}(\epsilon_{B})<-1.5$ and
$-4<\mbox{log}_{10}(\epsilon_{B})<-3.5$. The mean of this distribution is
$6.3\times 10^{-2}$ and the median is $1.0\times 10^{-2}$. The important point
of the $\epsilon_{B}$ histogram is that $\epsilon_{B}$ varies over 4 orders of
magnitude, showing that $\epsilon_{B}$ has a wide distribution and is an
uncertain parameter.
$\begin{array}[]{c}\begin{array}[]{cc}\includegraphics[scale={0.6}]{epsilon_e_histogram.eps}&\includegraphics[scale={0.6}]{epsilon_B_histogram.eps}\end{array}\end{array}$
Figure 1.— Left Panel: Histogram of the distribution of $\epsilon_{e}$ values
we found in the literature (Table LABEL:epsilon_e_and_epsilon_B_table in
Appendix A). Right Panel: Histogram of the distribution of $\epsilon_{B}$
values we found in the literature (Table LABEL:epsilon_e_and_epsilon_B_table
in Appendix A)
## 3\. Upper Limit On $\epsilon_{B}$ With Swift X-ray Light Curves
### 3.1. Constraining $\epsilon_{B}$ With The X-ray Light Curve Steep Decline
One interesting property found by Swift (Gehrels et al., 2004) is that at
early times, about 50% of the light curves detected by the XRT (X-ray
Telescope, Burrows et al., 2005) display a very rapid decline in flux, known
as the steep decline (Gehrels et al., 2009). The flux during the steep decline
typically decays as $t^{-3}$ and it usually lasts $\sim 10^{2}-10^{3}$ sec. By
extrapolating the BAT (Burst Alert Telescope, Barthelmy et al., 2005) emission
to the X-ray band, O’Brien et al. 2006 showed that there is a continuous
transition between the end of the prompt emission and the start of the steep
decline phase. This important conclusion lead to the interpretation that the
X-ray steep decline has an origin associated with the end of the prompt
emission.
The favored interpretation for the origin of the steep decline is high
latitude emission (Kumar & Panaitescu, 2000). Although high latitude emission
is able to explain most of the steep decline observations, some GRBs display
spectral evolution during the steep decline (Zhang et al., 2007b), which is
not expected. In any case, the steep decline cannot be produced by the
external-forward shock. Therefore, the observed flux during the steep decline
should be larger than or equal to the flux produced by the external-forward
shock. We do however assume that the time at which the steep decline typically
ends, at about $10^{2}-10^{3}$ seconds, is past the deceleration time222The
deceleration time marks the time when about half of the kinetic energy of the
blastwave has been transferred to the surrounding medium.. For our upper limit
on $\epsilon_{B}$ with X-ray data, we will use the expression for the flux
from the external-forward shock, which uses the kinetic energy of the
blastwave given by the Blandford and McKee solution (Blandford & McKee, 1976).
Since this solution is only valid for a decelerating blastwave, we need to be
past the deceleration time for it to be applicable.
Theoretically, it is expected that that the deceleration time occurs before
the end of of the steep decline. Depending on the density profile of the
surrounding medium, the deceleration time $t_{\mathrm{dec}}$ is given by
$t_{\mathrm{dec}}=\left\\{\begin{array}[]{ll}\hskip-7.0pt(220\mbox{
sec})E_{53}^{1/3}n_{0}^{-1/3}\Gamma_{2}^{-8/3}(1+z)&s=0\\\ &\\\
\hskip-7.0pt(67\mbox{
sec})E_{53}A_{*,-1}^{-1}\Gamma_{2}^{-4}(1+z)&s=2\end{array}\right.$ (2)
(e.g. Panaitescu & Kumar, 2000). In these expressions, $\Gamma$ is the Lorentz
factor of the shocked fluid, $z$ is the redshift, and we have adopted the
usual notation $Q_{n}\equiv Q/10^{n}$. For $s=2$, the proportionality constant
of the density, $A$, is normalized to the typical mass loss rate and stellar
wind velocity of a Wolf-Rayet star, which is denoted by $A_{*}$ and is defined
as $A_{*}\equiv A/(5\times 10^{11}\mbox{ g}\mbox{ cm}^{-1})$ (Chevalier & Li,
2000). For typical GRB afterglow parameters of $E_{53}=1$, $n_{0}=1$ (or
$A_{*}=0.1$ for $s=2$), $\Gamma_{2}=3$, and $z=2.5$, the deceleration time is
under $100$ seconds for both $s=0$ and $s=2$. Although there can be a large
uncertainty in the afterglow parameters $E$ and $n$, $\Gamma$ is the most
important parameter when calculating $t_{\mathrm{dec}}$ since
$t_{\mathrm{dec}}$ has a very strong dependence on $\Gamma$. For $s=0$
($s=2$), even if we take extreme parameters for $E$ and $n$, such as a high
kinetic energy of $E_{53}=100$ and a low density of $n=10^{-3}\mbox{ cm}^{-3}$
($A_{*}=10^{-2}$), with a typical $\Gamma_{2}=3-4$ (e.g. Molinari et al.,
2007; Xue et al., 2009; Liang et al., 2010), $t_{\mathrm{dec}}$ is still a few
hundred seconds. Thus, since the deceleration time is less than the typical
time at which the steep decline ends, the onset of the external-forward shock
emission occurs before the end of the steep decline.
Observationally, the deceleration time is also seen to occur before the end of
the steep decline for many GRBs. If the dominant contribution to the light
curves at early times is the external-forward shock, the light curve is
expected to rise as a power law, reach a peak, and then decline as a power
law, with the peak signifying the deceleration time. In Liang et al. 2010, a
sample of optical light curves that display this peak is studied. In their
Figure 1, for each GRB, they display both the optical light curve and the
X-ray light curve. For all their bursts that display a steep decline in the
X-ray light curve (except for GRB 080303 and GRB 081203A), it can be seen that
the peak of the optical light curve occurs before the end of the X-ray steep
decline. We did not include GRB 080330 and GRB 081203A in our samples.
Since an increase in $\epsilon_{B}$ increases $f_{\mathrm{ES}}$ (external-
forward shock flux), the condition that the X-ray flux during the steep
decline should be larger than or equal to $f_{\mathrm{ES}}$ gives an upper
limit on $\epsilon_{B}$. Since our goal is to attain the most stringent upper
limit on $\epsilon_{B}$, we take this constraint at the end of the steep
decline. Explicitly, the constraint we will use to find an upper limit on
$\epsilon_{B}$ with X-ray data is
$f_{\mathrm{EoSD}}\geq f_{\mathrm{ES}}(E,n,s,\epsilon_{e},\epsilon_{B},p).$
(3)
In this inequality, $f_{\mathrm{EoSD}}$ represents the observed flux at the
end of the steep decline (EoSD). We have also explicitly shown the dependence
of $f_{\mathrm{ES}}$ on the afterglow parameters. We will now discuss the
assumptions we make on the other afterglow parameters, which will allow us to
calculate an upper limit on $\epsilon_{B}$.
### 3.2. The Other Afterglow Parameters
#### 3.2.1 $E$ and $\epsilon_{e}$
Although we do not know the isotropic equivalent kinetic energy of the
blastwave $E$, we can calculate the isotropic energy released in gamma-rays
during the prompt emission:
$E_{iso}^{\gamma}=\frac{\mbox{fluence}\times 4\pi d_{L}^{2}}{1+z}.$ (4)
In this equation, the fluence has units of ergs/$\mbox{cm}^{2}$ and it
represents the flux detected in gamma-rays, integrated over the duration of
the prompt emission. $d_{L}$ is the luminosity distance. Since we are
interested in the fluence radiated in gamma-rays, we use the fluence detected
in the BAT band, ranging from 15-150 keV (Barthelmy et al., 2005). The
fluences detected by BAT for each GRB can be found in NASA’s Swift GRB Table
and Lookup website333
http://heasarc.gsfc.nasa.gov/docs/swift/archive/grb_table/ .
To convert $E_{iso}^{\gamma}$ to $E$, we need to know the efficiency in the
conversion of kinetic energy of the jet to prompt gamma-ray emission. Recent
studies on the prompt emission efficiency, using X-ray light curves with
plateaus detected by Swift, were presented in Granot et al. (2006) and Zhang
et al. (2007a). For the 23 GRBs for which Zhang et al. (2007a) presented
results for the efficiency (see their Table 3), more than half of them were
found to have a high efficiency $\gtrsim$ 30% 444In Table 3 of Zhang et al.
(2007a), they present two different estimates for the efficiency. If the
shallow decay of plateaus seen in X-ray light curves is due to energy
injection, the more appropriate of the two estimates for the efficiency is
denoted as $\eta_{\gamma}(t_{\mathrm{dec}})$.
$\eta_{\gamma}(t_{\mathrm{dec}})$ represents the efficiency in gamma-ray
radiation ($\eta_{\gamma}$) calculated at the deceleration time
$t_{\mathrm{dec}}$., with a few being estimated to have an efficiency as high
as 90%. A high efficiency of $\sim$90% was also found in Granot et al. (2006).
However, Fan & Piran 2006 argue that bursts with X-ray plateaus should have
more moderate efficiencies $\sim 10\%$. In addition, Zhang et al. (2007a)
mention that the efficiencies they calculate for some bursts have large errors
due to the uncertainty in the microphysical parameters $\epsilon_{e}$ and
$\epsilon_{B}$. In this work, we will take a standard choice and calculate $E$
with the expression
$E=5E_{\mathrm{iso}}^{\gamma}.$ (5)
From the definition of the efficiency,
$\eta=E^{\gamma}_{\mathrm{iso}}/(E^{\gamma}_{\mathrm{iso}}+E)$, Equation 5
corresponds to an efficiency of $\sim 20$%. At the end of Section 3.5, we will
discuss in more detail how the uncertainty in the efficiency affects our
results. For our X-ray sample (see Section 3.3), from Equation 5, we found
values for $E$ in the range of $10^{51}-10^{54}$ ergs, with a typical value
$\sim 10^{53}$ ergs. The average value of $E_{53}$ for our X-ray sample is 2.8
and the median is 1.6. The fluence detected in the BAT band,
$E^{\gamma}_{\mathrm{iso,}52}$, $z$, and $d_{L28}$ for each GRB in our X-ray
sample are shown in Table LABEL:x-ray_properties_table 555Except for GRBs
060708, 060906, 061021, 061222A, 080906, and 081230, we obtained all the
redshifts from NASA’s Swift GRB Table and Lookup website. For the exceptions,
we obtained the redshifts from the website on GRB redshifts maintained by J.
Greiner, http://www.mpe.mpg.de/ jcg/grbgen.html
http://www.mpe.mpg.de/$\sim$jcg/grbgen.html.
For $\epsilon_{e}$, we assumed a value of 0.2 for all of the GRBs in our
sample. This choice for $\epsilon_{e}$ is justified from the results of
$\epsilon_{e}$ with previous afterglow studies (Figure 1) and with recent
results from simulations (Sironi & Spitkovsky, 2011), as discussed in Section
2.
#### 3.2.2 Electron power-law index and density profile
In afterglow studies, $s$ and $p$ can be obtained by determining which closure
relation the observed temporal and spectral decay indices satisfy. We cannot
use this strategy to determine $p$ and $s$ for our X-ray sample since the
external-forward shock flux is below the observed steep decline emission.
Instead, we use a fixed $p$ for all GRBs in our X-ray sample. We consider a
small value of $p=2.2$, a typical value of $p=2.4$ for Swift GRBs (Curran et
al., 2010), and a larger value of $p=2.8$. Previous afterglow studies have
found that the majority of afterglow observations are better described by a
constant density medium (e.g. Panaitescu & Kumar, 2002; Schulze et al., 2011).
However, there are still a number of cases where the wind medium is a better
model for the afterglow observations. Therefore, we will consider both $s=0$
and $s=2$ when displaying the results for the upper limit on $\epsilon_{B}$.
#### 3.2.3 Density
The density of the medium in the vicinity of GRBs is a highly uncertain
parameter. A histogram of values of $n$ determined by previous afterglow
modelling studies can be found in Figure 9 of Soderberg et al. (2006), which
shows $n$ to vary over 5 orders of magnitude, ranging from $\sim 10^{-3}\mbox{
cm}^{-3}$ to $\sim 10^{2}\mbox{ cm}^{-3}$. In Section 3.5, we will discuss in
more detail how the uncertainty in the density affects our results for the
$\epsilon_{B}$ upper limits.
### 3.3. The X-ray Sample
For our constraint on $\epsilon_{B}$ with X-ray data, we only consider X-ray
data detected by the XRT on board Swift. We used the X-ray light curves
presented in Butler & Kocevski (2007)666 http://butler.lab.asu.edu/swift/
http://butler.lab.asu.edu/swift/. With the exception of two cases, we only
consider bursts that display a steep decline in their X-ray light curve (see
below). After the end of the steep decline, GRBs display a variety of temporal
decays (Evans et al., 2009). Our sample can be divided into 4 different
subgroups, based on the temporal decay after the steep decline:
1. 1.
Steep Decline To Plateau: In this subgroup, GRBs display a plateau after the
steep decline ($\sim$ 73% of our sample). In Table
LABEL:x-ray_properties_table, we display the time at the end of the steep
decline in units of $10^{2}$ sec, $t_{2,\mathrm{EoSD}}$, and the observed flux
at the end of the steep decline at 1 keV, $f_{1\mathrm{keV},\mathrm{EoSD}}$,
in $\mu$Jy.
2. 2.
Steep Decline To Normal Decline: In this subgroup, GRBs display a temporal
decay of $\alpha\sim 1$ after the steep decline ($\sim$ 18% of our sample). In
Table LABEL:x-ray_properties_table, we show the time and the flux at the end
of the steep decline.
3. 3.
Clear Steep Decline But Not A Clear End To The Steep Decline: In this
subgroup, it cannot be determined where the steep decline ends (the XRT
observations end before the steep decline ends). The following GRBs fall into
this subgroup: 050315, 060202, 070419A, 071122, and 090516. For these GRBs, in
Table LABEL:x-ray_properties_table, we show the flux and the time of the last
steep decline point observed. For these GRBs, to be sure that the last steep
decline point observed is past the deceleration time, we made sure that it is
at a few hundred seconds.
4. 4.
No Clear Steep Decline Seen, Just Plateau: Two of the GRBs in our X-ray
sample, 050401 and 060927, do not display a steep decline. The first
observation of the X-ray light curve for these bursts is during the plateau.
We did not remove these GRBs from our X-ray sample because they are also part
of our optical sample (bursts with both X-ray and optical data are important
because they allow us to cross-check our results, see Section 5). For these 2
bursts, we considered the first observation in the X-ray light curve for our
constraint so that we have the least amount of energy injection. We made sure
that this point is at least at a few hundred seconds so that we can be
confident that the onset of the external-forward shock emission has occurred.
For these two bursts, we show the time and the flux of the first X-ray
observation of the plateau in Table LABEL:x-ray_properties_table.
In addition, 25% of the GRBs in our X-ray sample display X-ray flares during
the steep decline. We only consider bursts where the X-ray flare ends before
the end of the steep decline because it is difficult to determine the flux and
time at the end of the steep decline for bursts that show X-ray flares near
the end of the steep decline. It is fine to consider these bursts because
observationally, after the X-ray flare, the X-ray light curve is seen to
return to the same temporal decay prior to the X-ray flare (e.g. Chincarini et
al. 2007). Lastly, GRB 051221A is the only short GRB in our X-ray sample; all
the other bursts in our X-ray sample are long GRBs.
Table 1Properties of X-ray Sample GRB | $z$ | $d_{L28}$ | Fluence | $E^{\gamma}_{\mathrm{iso,}52}$ | $t_{2,\mathrm{EoSD}}$ | $f_{1\mathrm{keV,EoSD}}$ | $\mbox{log}_{10}(\epsilon_{B})$ | $\mbox{log}_{10}(\epsilon_{B})$
---|---|---|---|---|---|---|---|---
| | | $[\times 10^{-6}\mbox{ ergs}/\mbox{cm}^{2}$] | | | [$\mu$Jy] | ($s=0$) | ($s=2$)
050315 | 1.949 | 4.71 | 3.22 | 3.04 | 4 | 3 | -5.2 | -6.0
050401 | 2.9 | 7.67 | 8.22 | 15.56 | 2 | 80 | -4.7 | -5.4
050721 | 2.5 | 6.40 | 3.62 | 5.32 | 4 | 30 | -4.2 | -5.0
050803 | 0.422 | 0.71 | 2.15 | 0.10 | 3 | 10 | -3.8 | -5.4
050814 | 5.3 | 15.76 | 2.01 | 9.95 | 9 | 1 | -5.4 | -6.0
051221A | 0.547 | 0.97 | 1.15 | 0.09 | 6 | 5 | -3.5 | -4.9
060108 | 2.03 | 4.95 | 0.37 | 0.38 | 6 | 0.8 | -4.2 | -5.5
060111B | 2 | 4.86 | 1.60 | 1.58 | 2 | 8 | -4.6 | -5.8
060115 | 3.53 | 9.72 | 1.71 | 4.48 | 8 | 1 | -5.2 | -5.9
060210 | 3.91 | 10.99 | 7.66 | 23.66 | 8 | 40 | -4.4 | -4.7
060418 | 1.49 | 3.37 | 8.33 | 4.79 | 5 | 40 | -4.3 | -4.9
060502A | 1.51 | 3.43 | 2.31 | 1.36 | 3 | 7 | -4.6 | -5.7
060607A | 3.082 | 8.25 | 2.55 | 5.35 | 5 | 60 | -3.6 | -4.3
060707 | 3.425 | 9.37 | 1.60 | 3.99 | 9 | 3 | -4.5 | -5.2
060708 | 1.92 | 4.62 | 0.49 | 0.45 | 2 | 10 | -3.7 | -5.2
060714 | 2.71 | 7.06 | 2.83 | 4.78 | 3 | 10 | -4.8 | -5.7
060729 | 0.54 | 0.96 | 2.61 | 0.20 | 6 | 7 | -3.8 | -5.1
060814 | 0.84 | 1.65 | 14.60 | 2.71 | 8 | 8 | -5.0 | -5.6
060906 | 3.686 | 10.24 | 2.21 | 6.21 | 4 | 0.8 | -5.9 | -6.7
060926 | 3.208 | 8.66 | 0.22 | 0.49 | 2 | 4 | -3.8 | -5.4
060927 | 5.6 | 16.81 | 1.13 | 6.08 | 0.8 | 8 | -5.3 | -6.6
061021 | 0.3463 | 0.56 | 2.96 | 0.09 | 3 | 10 | -3.9 | -5.5
061110A | 0.758 | 1.45 | 1.06 | 0.16 | 5 | 5 | -3.6 | -6.0
061121 | 1.314 | 2.88 | 13.70 | 6.19 | 2 | 40 | -5.1 | -5.8
061222A | 2.088 | 5.13 | 7.99 | 8.55 | 2 | 40 | -4.9 | -5.7
070110 | 2.352 | 5.94 | 1.62 | 2.14 | 4 | 8 | -4.3 | -5.3
070306 | 1.497 | 3.39 | 5.38 | 3.12 | 7 | 4 | -5.0 | -5.6
070714B | 0.92 | 1.85 | 0.72 | 0.16 | 4 | 6 | -3.5 | -5.0
070802 | 2.45 | 6.24 | 0.28 | 0.40 | 5 | 0.8 | -4.2 | -5.5
071122 | 1.14 | 2.41 | 0.58 | 0.20 | 8 | 0.8 | -4.1 | -5.4
080310 | 2.43 | 6.18 | 2.30 | 3.22 | 10 | 2 | -4.8 | -5.4
080413A | 2.433 | 6.19 | 3.50 | 4.91 | 2 | 10 | -5.1 | -6.0
080430 | 0.767 | 1.47 | 1.20 | 0.19 | 2 | 10 | -3.9 | -5.5
080607 | 3.036 | 8.10 | 24.00 | 49.07 | 3 | 90 | -5.2 | -5.5
080721 | 2.591 | 6.68 | 12.00 | 18.75 | 0.7 | 900 | -4.3 | -5.1
080905B | 2.374 | 6.00 | 1.80 | 2.42 | 2 | 20 | -4.3 | -5.4
080906 | 2.1 | 5.16 | 3.50 | 3.78 | 7 | 20 | -4.0 | -4.7
080916A | 0.689 | 1.29 | 4.00 | 0.50 | 3 | 10 | -4.4 | -5.7
081007 | 0.5295 | 0.93 | 0.71 | 0.05 | 2 | 8 | -3.5 | -5.3
081008 | 1.9685 | 4.77 | 4.30 | 4.14 | 6 | 20 | -4.2 | -4.9
081230 | 2 | 4.86 | 0.82 | 0.81 | 3 | 6 | -4.1 | -5.3
090418A | 1.608 | 3.71 | 4.60 | 3.05 | 2 | 20 | -4.8 | -5.8
090516A | 4.109 | 11.66 | 9.00 | 30.08 | 6 | 10 | -5.4 | -5.7
090519 | 3.85 | 10.79 | 1.20 | 3.62 | 5 | 0.8 | -5.4 | -6.2
090529 | 2.625 | 6.79 | 0.68 | 1.09 | 20 | 0.5 | -4.3 | -5.1
090618 | 0.54 | 0.96 | 105.00 | 7.86 | 4 | 200 | -4.9 | -5.3
090926B | 1.24 | 2.68 | 7.30 | 2.95 | 5 | 10 | -4.8 | -5.6
091029 | 2.752 | 7.19 | 2.40 | 4.16 | 6 | 1 | -5.5 | -6.2
091109A | 3.076 | 8.25 | 1.60 | 3.35 | 5 | 2 | -5.0 | -5.8
100302A | 4.813 | 14.06 | 0.31 | 1.33 | 8 | 1 | -4.2 | -5.2
100425A | 1.755 | 4.14 | 0.47 | 0.37 | 3 | 7 | -3.6 | -5.0
100513A | 4.772 | 13.92 | 1.40 | 5.91 | 9 | 7 | -4.1 | -4.8
100621A | 0.542 | 0.96 | 21.00 | 1.58 | 4 | 20 | -5.0 | -5.8
100704A | 3.6 | 9.95 | 6.00 | 16.23 | 6 | 9 | -5.1 | -5.5
100814A | 1.44 | 3.23 | 9.00 | 4.85 | 6 | 9 | -5.0 | -5.6
100906A | 1.727 | 4.05 | 12.00 | 9.09 | 3 | 20 | -5.2 | -5.9
110808A | 1.348 | 2.98 | 0.33 | 0.16 | 3 | 3 | -3.6 | -5.2
110818A | 3.36 | 9.16 | 4.00 | 9.67 | 20 | 1 | -5.3 | -5.5
111008A | 4.9898 | 14.68 | 5.30 | 23.95 | 3 | 9 | -5.5 | -6.1
111228A | 0.72 | 1.36 | 8.50 | 1.15 | 5 | 8 | -4.8 | -5.7
Table 1Properties of X-ray Sample (Continued)
### 3.4. Expected External-Forward Shock Emission At The End Of The Steep
Decline
The synchrotron afterglow spectrum consists of four power-law segments that
are smoothly joined together at three characteristic frequencies of
synchrotron emission (e.g. Sari, Piran, & Narayan, 1998; Granot & Sari, 2002).
These three characteristic frequencies are: $\nu_{a}$, the synchrotron self-
absorption frequency, $\nu_{i}$ (also commonly referred to as $\nu_{m}$), the
frequency of the photons emitted by the power-law distribution of injected
electrons with the smallest energy, and $\nu_{c}$, the cooling frequency
corresponding to electrons cooling on a dynamical time. For this work, we will
only consider the standard case for the ordering of the characteristic
frequencies, the slow cooling case, where $\nu_{a}<\nu_{i}<\nu_{c}$. One
argument against the fast cooling case ($\nu_{a}<\nu_{c}<\nu_{i}$) is that if
the observing frequency is between $\nu_{c}$ and $\nu_{i}$, the spectrum
should be $f_{\nu}\propto\nu^{-1/2}$; however, the spectral index $\beta=0.5$
disagrees with the typical observed afterglow spectral index $\beta\approx
0.9$ (e.g. Piro, 2001). The next question we need to consider is where the
X-ray band lies at the end of the steep decline (here, we consider 1 keV for
the X-ray band because the light curves we used are plotted at this energy).
The two possibilities for the spectral regime of the X-ray band are
$\nu_{i}<1\mbox{ keV}<\nu_{c}$ or $\nu_{c}<1\mbox{ keV}$. We rule out
$\nu_{c}<1\mbox{ keV}$ with the following two arguments.
First, we compare the observed flux at the end of the steep decline,
$f_{\mathrm{1keV,EoSD}}$, to the flux predicted by the external-forward shock
at the same time, if $\nu_{c}<1\mbox{ keV}$ (defined as $f_{\mathrm{pred}}$).
For $s=0$ , $f_{\mathrm{pred}}$ is given by (Granot & Sari, 2002)
$\displaystyle
f_{\mathrm{pred}}=0.855(p-0.98)10^{\frac{3p+2}{4}}8.64^{\frac{3p-2}{4}}e^{1.95p}$
$\displaystyle\times(1+z)^{\frac{p+2}{4}}(\bar{\epsilon}_{e,-1})^{p-1}\epsilon_{B}^{\frac{p-2}{4}}E_{53}^{\frac{p+2}{4}}(t_{\mathrm{2,EoSD}})^{-\frac{(3p-2)}{4}}$
$\displaystyle\times d_{L28}^{-2}\nu_{14}^{-\frac{p}{2}}\mbox{ mJy ,}$ (6)
where $\bar{\epsilon}_{e,-1}\equiv(p-2)/(p-1)\epsilon_{e,-1}$. For
$\nu_{c}<1\mbox{ keV}$, the external-forward shock flux is independent of the
density and the $s=2$ expression is almost identical. When calculating
$f_{\mathrm{pred}}$, for each of the bursts in our X-ray sample, we assumed a
standard $p=2.4$, $\epsilon_{e}=0.2$, and $\nu_{14}=2.4\times 10^{3}$, the
frequency corresponding to 1 keV. For the parameters $E$, $t$, $z$, and
$d_{L}$, we used the values given in Table LABEL:x-ray_properties_table for
each burst ($E=5E_{iso}^{\gamma}$). The remaining parameter we need to compute
$f_{\mathrm{pred}}$ is $\epsilon_{B}$. Since $\epsilon_{B}$ is raised to the
power of $(p-2)/4$, for a typical $p\sim 2-3$, the dependence on
$\epsilon_{B}$ is very weak. When calculating $f_{\mathrm{pred}}$, we assumed
a low value of $\epsilon_{B}=10^{-3}$.
We computed the ratio $f_{\mathrm{pred}}/f_{\mathrm{1keV,EoSD}}$ for all the
GRBs in our X-ray sample and found that
$f_{\mathrm{pred}}/f_{\mathrm{1keV,EoSD}}>1$ for all the bursts and
$f_{\mathrm{pred}}/f_{\mathrm{1keV,EoSD}}>10$ for $54/60$ bursts. The mean
value of $f_{\mathrm{pred}}/f_{\mathrm{1keV,EoSD}}$ is 50 and the median value
is 34. This means that the predicted flux from the external-forward shock,
when $\nu_{c}<\mbox{1 keV}$, over-predicts the observed flux at the end of the
steep decline by a factor that is larger than 10 for the majority of the
bursts. Therefore, the assumption that $\nu_{c}<1$ keV is incorrect. This is a
robust conclusion because the X-ray flux from the external-forward shock, when
$\nu_{c}<1$ keV, basically only depends on $\sim\epsilon_{e}E$ (see Equation
6), which cannot be decreased by a factor of $>10$ without introducing serious
efficiency problems in producing the prompt gamma-rays. Even if we allow for
an uncertainty of a factor of $\sim 2-3$ in both $\epsilon_{e}$ and $E$, this
is not enough to decrease $f_{\mathrm{pred}}/f_{\mathrm{1keV,EoSD}}$ below 1
for the majority of bursts in our X-ray sample.
Before continuing, we want to add that $f_{\mathrm{pred}}$ (Equation 6) also
has a dependence on the Compton-$Y$ parameter:
$f_{\mathrm{pred}}\propto(1+Y)^{-1}$. With this dependence, if the Compton-$Y$
parameter is large, then it is possible for $f_{\mathrm{pred}}$ to decrease
below $f_{\mathrm{1keV,EoSD}}$. For a few bursts in our X-ray sample, we
performed a detailed numerical calculation of the external-forward shock flux
with the formalism presented in Barniol Duran & Kumar (2011), which includes a
detailed treatment of Compton-$Y$ with Klein-Nishina effects. From this
calculation, we also found that $f_{\mathrm{pred}}$ over-predicts
$f_{\mathrm{1keV,EoSD}}$ by a factor larger than 10, which means that the
Inverse Compton (IC) cooling of electrons producing 1 keV synchrotron photons
is a weak effect (even when $\epsilon_{B}$ is small). Numerically, we also
found that without making any assumption about the location of $\nu_{c}$,
solutions to the constraint $f_{\mathrm{ES,1keV}}\leq f_{1keV,EoSD}$ were only
found when 1 keV $<\nu_{c}$ (when $\epsilon_{B}$ is small such that 1 keV
$<\nu_{c}$, it turns out that IC cooling of electrons producing 1 keV photons
takes place in the Klein-Nishina regime at these early times of $\sim$ few
$\times 100$ sec).
Another argument in favor of the spectral regime of the X-ray band being
$\nu_{i}<\mbox{1 keV}<\nu_{c}$ at the end of the steep decline comes from the
extrapolation of $\nu_{c}$ at late times to the end of the steep decline. In
Liang et al. (2008), they made fits to late time XRT light curves during the
normal decline phase and they also provided the value of $\nu_{c}$ during the
geometrical midpoint of the normal decline phase777The geometrical midpoint of
the normal decline is defined by
$\mbox{log}_{10}t=(\mbox{log}_{10}t_{1}+\mbox{log}_{10}t_{2})/2$ , where
$t_{1}$ represents the time of the first observation of the normal decline and
$t_{2}$ represents the time of the last observation of the normal decline..
For this argument, we focus on GRBs that are in common to our X-ray sample and
to the sample of Liang et al. (2008). For these bursts, we extrapolate
$\nu_{c}$ at late times to the end of the steep decline. In Liang et al.
(2008), they only considered a constant density medium, where $\nu_{c}\propto
t^{-1/2}$. The results of the extrapolation of $\nu_{c}$ are shown in Table 2.
GRB | $t_{\mathrm{late},4}$ | $t_{\mathrm{EoSD},2}$ | $\nu_{\mathrm{c,late}}$ (keV) | $\nu_{\mathrm{c,EoSD}}$ (keV)
---|---|---|---|---
050315 | $3.5$ | 4 | 0.17 | 1.6
050401 | $6.9$ | 2 | 4.1 | 77
051221A | $1.7$ | 6 | 4.1 | 22
060210 | $3.2$ | 8 | 0.97 | 6.1
060502A | $0.42$ | 3 | 4.1 | 15
060714 | $3.5$ | 3 | 3.6 | 39
060729 | $40$ | 6 | 1.4 | 36
060814 | $2.0$ | 8 | 4.1 | 21
060926 | $0.26$ | 2 | 6.7 | 24
061121 | $1.1$ | 2 | 4.2 | 31
070110 | $14$ | 4 | 4 | 74
Table 2This table shows the extrapolation of $\nu_{c}$ from late times to the
end of the steep decline for GRBs in common to our sample and to the sample in
Liang et al. (2008) (first column). The second column shows
$t_{\mathrm{late},4}$, the late time in units of $10^{4}$ seconds at which
Liang et al. (2008) determined $\nu_{c}$. The third column shows
$t_{\mathrm{EoSD},2}$, the time at the end of the steep decline in units of
$10^{2}$ seconds. $\nu_{\mathrm{c,late}}$, given in keV, is the value found in
Liang et al. (2008) for $\nu_{c}$ at $t_{\mathrm{late}}$. The last column
shows $\nu_{\mathrm{c,EoSD}}$ in keV. $\nu_{\mathrm{c,EoSD}}$ is found by
extrapolating $\nu_{\mathrm{c,late}}$ to $t_{\mathrm{EoSD}}$. Since Liang et
al. (2008) assume a constant density medium, we take a constant density medium
for all the GRBs in this sample when making the extrapolation of $\nu_{c}$ to
$t_{\mathrm{EoSD}}$.
In Table 2, we find that at the end of the steep decline, 1 keV
$<\nu_{c,\mathrm{EoSD}}$ for all GRBs. This further confirms our choice that
$\nu_{i}<\mbox{ 1 keV}<\nu_{c}$ at the end of the steep decline888Including
energy injection will make the values of $\nu_{\mathrm{c,EoSD}}$ in Table 2
larger, making the conclusion that $\nu_{i}<\mbox{ 1 keV}<\nu_{c}$ at the end
of the steep decline more robust. For the bursts that have plateaus in their
X-ray light curve, energy injection needs to be considered. During the energy
injection episode, $E$ increases as $E\propto t^{1-q}$ (Zhang et al., 2006),
where $q$ is a positive constant that satisfies $0\leq q\leq 1$. Therefore,
between the end of the plateau and the end of the steep decline, since
$\nu_{c}\propto E^{-1/2}t^{-1/2}$, $\nu_{c}\propto t^{-1+(q/2)}$. This time
evolution of $\nu_{c}$ is steeper than $\nu_{c}\propto t^{-1/2}$ without
energy injection for $s=0$ (Liang et al. (2008) only considered $s=0)$..
The knowledge of the spectral regime regime at the end of the steep decline
allows us to write in an explicit expression for $f_{\mathrm{ES}}$ in Equation
3:
$\frac{f_{\mathrm{1keV,EoSD}}}{\mbox{mJy}}\geq\left\\{\begin{array}[]{ll}\hskip-7.0pt0.461(p-0.04)10^{\frac{3p+1}{4}}8.64^{\frac{3(p-1)}{4}}e^{2.53p}&\\\
\times(1+z)^{\frac{p+3}{4}}(\bar{\epsilon}_{e,-1})^{p-1}\epsilon_{B}^{\frac{p+1}{4}}n_{0}^{\frac{1}{2}}E_{53}^{\frac{p+3}{4}}&\\\
\times(t_{\rm
2,EoSD})^{-\frac{3(p-1)}{4}}d_{L28}^{-2}\nu_{14}^{-\frac{(p-1)}{2}}\mbox{
$s=0$}\\\ &\\\
\hskip-7.0pt3.82(p-0.18)10^{\frac{3p-1}{4}}8.64^{\frac{3p-1}{4}}e^{2.54p}&\\\
\times(1+z)^{\frac{p+5}{4}}(\bar{\epsilon}_{e,-1})^{p-1}\epsilon_{B}^{\frac{p+1}{4}}A_{*,-1}E_{53}^{(p+1)/4}&\\\
\times(t_{\rm
2,EoSD})^{-\frac{(3p-1)}{4}}d_{L28}^{-2}\nu_{14}^{-\frac{(p-1)}{2}}\mbox{
$s=2$}.\end{array}\right.$ (7)
On the left hand side of this inequality, we have the observed X-ray flux at
the end of the steep decline and on the right hand side we have the expression
for the external-forward shock flux when $\nu_{i}<\mbox{ 1 keV}<\nu_{c}$
(Granot & Sari, 2002). The notation used for $\bar{\epsilon}_{e}$ and $A_{*}$
is defined as $\bar{\epsilon}_{e,-1}\equiv\bar{\epsilon}_{e}/10^{-1}$ and
$A_{*,-1}\equiv A_{*}/10^{-1}$. The expressions in Equation 7 are only valid
for $p>2$, which we consider in this work (for $p<2$, see Bhattacharya 2001
and Resmi & Bhattacharya 2008).
Before displaying our results for the $\epsilon_{B}$ upper limit for our
entire X-ray sample, we show a simple calculation to get an idea of what
values to expect for the $\epsilon_{B}$ upper limits from the X-ray constraint
given in Equation 7. For a standard $p=2.4$, the X-ray constraint is
$\frac{f_{\mathrm{1keV,EoSD}}}{\mbox{mJy}}\geq\left\\{\begin{array}[]{ll}\hskip-7.0pt5.0\times
10^{1}t_{2,\mathrm{EoSD}}^{-1.05}\epsilon_{e,-1}^{1.4}E_{53}^{1.35}n_{0}^{0.5}\epsilon_{B}^{0.85}&s=0\\\
&\\\ \hskip-7.0pt7.0\times
10^{2}t_{2,\mathrm{EoSD}}^{-1.55}\epsilon_{e,-1}^{1.4}E_{53}^{0.85}A_{*,-1}\epsilon_{B}^{0.85}&s=2.\end{array}\right.$
(8)
For this calculation, we used the average $z=2.5$ for Swift GRBs (Gehrels et
al., 2009) (with a corresponding $d_{L28}=6.4$) and $\nu_{14}$ corresponding
to 1 keV. Solving for $\epsilon_{B}$, the upper limit depends on the afterglow
parameters as999For bursts with plateaus in their X-ray light curve, it is
possible that energy injection begins before the steep decline ends. However,
even with energy injection, there still exists a self-similar solution for the
energy (Equation 52 of Blandford & McKee 1976). For both $s=0$ and $s=2$, we
calculated the external-forward shock synchrotron flux with this new self-
similar solution and then calculated the upper limit on $\epsilon_{B}$ as
shown in Equation 9. We found that the $\epsilon_{B}$ upper limits are
affected by less than a factor of $\sim 2$. Thus, even if energy injection
begins before the end of the steep decline, it has very little to no effect on
our $\epsilon_{B}$ upper limit results.
$\epsilon_{B}\leq\left\\{\begin{array}[]{ll}\hskip-7.0pt1.0\times
10^{-2}\left(\frac{f_{1keV,\mathrm{EoSD}}}{\mbox{mJy}}\right)^{1.18}t_{2,\mathrm{EoSD}}^{1.24}&\\\
\times\epsilon_{e,-1}^{-1.65}E_{53}^{-1.59}n_{0}^{-0.59}&s=0\\\ &\\\
\hskip-7.0pt4.5\times
10^{-4}\left(\frac{f_{1keV,\mathrm{EoSD}}}{\mbox{mJy}}\right)^{1.18}t_{2,\mathrm{EoSD}}^{1.82}&\\\
\times\epsilon_{e,-1}^{-1.65}E_{53}^{-1}A_{*,-1}^{-1.18}&s=2.\end{array}\right.$
(9)
The median values for our X-ray sample for the parameters
$f_{\mathrm{1keV,EoSD}}$, $t_{2,\mathrm{EoSD}}$, and $E_{53}$ are $8\times
10^{-3}$ mJy, 4, and 1.6, respectively. Using these median values and
$\epsilon_{e,-1}=2$, the upper limit on $\epsilon_{B}$ becomes
$\epsilon_{B}\leq\left\\{\begin{array}[]{ll}\hskip-7.0pt2.8\times
10^{-5}\times n_{0}^{-0.59}&s=0\\\ &\\\ \hskip-7.0pt3.7\times 10^{-6}\times
A_{*-1}^{-1.18}&s=2.\end{array}\right.$ (10)
For a standard $n_{0}=1$ and $A_{*,-1}=1$, it can be seen that the
$\epsilon_{B}$ upper limit is lower for $s=2$. This is expected because for
$A_{*,-1}=1$, there is a larger density for the surrounding medium within a
typical deceleration radius of $10^{17}\mbox{ cm}$.
In Equation 10, the explicit dependence of the $\epsilon_{B}$ upper limit on
the density is shown for $p=2.4$. In the next subsection, we will display the
results of the $\epsilon_{B}$ upper limits for our entire X-ray sample. To
keep the density-dependence, we will display histograms of upper limits on the
quantity $\epsilon_{B}n_{0}^{0.59}$ ($\epsilon_{B}A_{*-1}^{1.18}$) for $s=0$
($s=2$) for $p=2.4$, or $\epsilon_{B}n_{0}^{2/(p+1)}$
($\epsilon_{B}A_{*-1}^{4/(p+1)}$) for $s=0$ ($s=2$) for a general $p$ (see
Equation 7).
### 3.5. $\epsilon_{B}$ Upper Limits For Our X-ray Sample
We display the results for the upper limits (from Equation 7) on the quantity
$\epsilon_{B}n_{0}^{2/(p+1)}$ ($\epsilon_{B}A_{*-1}^{4/(p+1)}$) for $s=0$
($s=2$) assuming all GRBs in our X-ray sample have $p=2.2$, 2.4, and 2.8 in
the Top-Left, Top-Right, and Bottom panels of Figure 2, respectively. Two
histograms are shown in each panel, one for $s=0$ and the other for $s=2$.
Table 3 shows the mean and median upper limits on the quantity
$\epsilon_{B}n_{0}^{2/(p+1)}$ ($\epsilon_{B}A_{*-1}^{4/(p+1)}$) for $s=0$
($s=2$) for each histogram.
X-ray | | | |
---|---|---|---|---
$(s=0)$ | $p=2.2$ | $p=2.4$ | $p=2.8$ |
Mean | $1.1\times 10^{-4}$ | $7.2\times 10^{-5}$ | $6.1\times 10^{-5}$ |
Median | $4.2\times 10^{-5}$ | $2.8\times 10^{-5}$ | $2.0\times 10^{-5}$ |
X-ray | | | |
$(s=2)$ | $p=2.2$ | $p=2.4$ | $p=2.8$ |
Mean | $7.9\times 10^{-6}$ | $5.5\times 10^{-6}$ | $5.6\times 10^{-6}$ |
Median | $4.0\times 10^{-6}$ | $3.2\times 10^{-6}$ | $3.4\times 10^{-6}$ |
Opt. | | | |
$(s=0)$ | $p=2.2$ | $p=2.4$ | $p=2.8$ | $p$ from $\alpha_{\mathrm{O}}$
Mean | $3.5\times 10^{-4}$ | $1.0\times 10^{-4}$ | $2.9\times 10^{-5}$ | $9.5\times 10^{-5}$
Median | $1.1\times 10^{-4}$ | $3.3\times 10^{-5}$ | $5.5\times 10^{-6}$ | $2.4\times 10^{-5}$
Opt. | | | |
$(s=2)$ | $p=2.2$ | $p=2.4$ | $p=2.8$ |
Mean | $7.2\times 10^{-5}$ | $2.4\times 10^{-5}$ | $7.5\times 10^{-6}$ |
Median | $1.2\times 10^{-5}$ | $3.9\times 10^{-6}$ | $8.8\times 10^{-7}$ |
Table 3Mean and median $\epsilon_{B}$ values for the X-ray (upper limits on
$\epsilon_{B}$) and optical (measurements of $\epsilon_{B}$) histograms shown
in Figures 2 and 3. The section labeled “X-ray $(s=0)$” (“X-ray $(s=2)$”)
shows the mean and median $\epsilon_{B}$ upper limits assuming a constant
density (wind) medium with a standard $n=1\mbox{ cm}^{-3}$ ($A_{*}=0.1$). The
columns show the value of $p$ that was assumed. The section labeled “Opt.
$(s=0)$” (“Opt. $(s=2)$”) shows the mean and median $\epsilon_{B}$
measurements assuming a constant density (wind) medium with a standard
$n=1\mbox{ cm}^{-3}$ ($A_{*}=0.1$). The columns show the value of $p$ that was
assumed. The column labeled “$p$ from $\alpha_{O}$” shows the mean and median
$\epsilon_{B}$ measurements with $p$ determined from $\alpha_{O}$. There are
60 GRBs in our X-ray sample and 35 GRBs in our optical sample.
For the remainder of this section, we assume a standard $n_{0}=1$
($A_{*-1}=1$) for $s=0$ ($s=2$) when discussing our results for the
$\epsilon_{B}$ upper limits.
The $\epsilon_{B}$ upper limit histograms show a wide distribution. For a
constant density (wind) medium, all the histograms show a distribution ranging
from $\sim 10^{-6}-10^{-3}$ ($\sim 10^{-7}-10^{-4}$). For a constant density
(wind) medium, the mean and median $\epsilon_{B}$ upper limit values are
$\sim\mbox{ few}\times 10^{-5}$ ($\sim\mbox{ few}\times 10^{-6}$). Assuming a
different value of $p$ does not have a significant effect on the distributions
of the $\epsilon_{B}$ upper limits for our X-ray sample. For both the $s=0$
and $s=2$ cases, when changing $p$, the mean and median $\epsilon_{B}$ upper
limit values change by less than a factor $\sim 2$. Although previous
afterglow studies also showed a wide distribution for $\epsilon_{B}$ (Figure
1), our distribution of $\epsilon_{B}$ upper limits is shifted towards lower
values. Unlike Figure 1, which shows that many GRBs have been reported to have
$\epsilon_{B}\sim 10^{-3}-10^{-1}$, none of our histograms of $\epsilon_{B}$
upper limits show an $\epsilon_{B}$ upper limit larger than $10^{-3}$.
We now discuss how our assumptions on the afterglow parameters can affect the
distribution of $\epsilon_{B}$ upper limits. For this discussion, we will take
a typical $p=2.4$; Equation 9 shows how the $\epsilon_{B}$ upper limit depends
on the other afterglow parameters. The strongest dependence is on
$\epsilon_{e}$, which is raised to the power of $-1.65$. However, as we
displayed in Figure 1, according to previous studies, the distribution of
$\epsilon_{e}$ values is narrow, with the $\epsilon_{e}$ values ranging only
over one order of magnitude. In addition, $\sim 62\%$ of the bursts have
$\epsilon_{e}\sim 0.1-0.3$. From Figure 1, a likely error in $\epsilon_{e}$
from our assumed $\epsilon_{e}=0.2$ is a factor $\sim 2$. From Equation 9, an
error in $\epsilon_{e}$ by a factor $\sim 2$ will only lead to an error in the
$\epsilon_{B}$ upper limit by a factor $\sim 3$. For a constant density (wind)
medium, the $\epsilon_{B}$ upper limit depends on $E$ as $E^{-1.59}$
($E^{-1}$). We assumed an efficiency of $\sim 20\%$ in the conversion of
kinetic energy to prompt gamma-ray radiation. Recent studies have found higher
values for the efficiency (Granot et al. 2006, Zhang et al. 2007a, see
however, Fan & Piran 2006). In Zhang et al. 2007a, the mean (median)
efficiency they reported is $\sim 37\%$ ($\sim 32\%$). Taking the efficiency
to be $\sim 30\%-40\%$ instead of $\sim 20\%$ would lead to an error in $E$ by
a factor $\sim 2-3$. From Equation 9, an error in $E$ by a factor $\sim 2-3$
would lead to an error in the $\epsilon_{B}$ upper limit by a factor $\sim
3-6$ ($\sim 2-3$) for a constant density (wind) medium. Lastly, the largest
source of uncertainty for the $\epsilon_{B}$ upper limits is the density,
since it has been reported to have a range $\sim 10^{-3}\mbox{
cm}^{-3}-10^{2}\mbox{ cm}^{-3}$. For a constant density (wind) medium, the
$\epsilon_{B}$ upper limit depends on the density as $n^{-0.59}$
($A_{*}^{-1.18}$). For $s=0$ ($s=2$), we assumed a standard $n=1\mbox{
cm}^{-3}$ ($A_{*}=0.1$). An error in the density by a factor $\sim 10^{3}$
($\sim 10^{2}$) will lead to an error in the $\epsilon_{B}$ upper limit by a
factor $\sim 60$ ($\sim 230$).
In summary, the expected errors in $\epsilon_{e}$ and $E$ of a factor $\sim
2-3$ will not change the $\epsilon_{B}$ upper limits by an order of magnitude.
On the other hand, the density is a very uncertain parameter and an error in
the density by $\sim 2-3$ orders of magnitude will lead to an error in the
$\epsilon_{B}$ upper limits by $\sim 1-2$ orders of magnitude.
$\begin{array}[]{c}\begin{array}[]{cc}\includegraphics[scale={0.6}]{x_ray_p_2_2_histogram.eps}&\includegraphics[scale={0.6}]{x_ray_p_2_4_histogram.eps}\end{array}\\\
\includegraphics[scale={0.6}]{x_ray_p_2_8_histogram.eps}\end{array}$
Figure 2.— The Top-Left, Top-Right, and Bottom panels show the histograms of
upper limits assuming $p=2.2$, $p=2.4$, and $p=2.8$ respectively, for all of
the GRBs in our X-ray sample (obtained with Equation 7). The filled-in (un-
filled) histograms show upper limits on the quantity
$\epsilon_{B}n_{0}^{2/(p+1)}$ ($\epsilon_{B}A_{*-1}^{4/(p+1)}$) for $s=0$
($s=2$) assuming all the GRBs in our X-ray sample are described by a constant
density (wind) medium.
Two additional parameters that can affect our $\epsilon_{B}$ upper limits are:
1. $\xi$, the fraction of electrons accelerated to a power-law distribution
and 2. $f$, which is a factor that takes into account the degeneracy for a set
of afterglow parameters. For a set of parameters $E$, $n$, $\epsilon_{e}$,
$\epsilon_{B}$, $\xi$ producing the observed external-forward shock flux,
another set of primed parameters $E^{\prime}=E/f$, $n^{\prime}=n/f$,
$\epsilon_{e}^{\prime}=f\epsilon_{e}$, $\epsilon_{B}^{\prime}=f\epsilon_{B}$,
$\xi^{\prime}=f\xi$ can also produce the observed external-forward shock flux
(Eichler & Waxman, 2005). Afterglow studies usually assume $\xi=f=1$ for
simplicity but $\xi\leq 1$ and $m_{e}/m_{p}\leq f\leq 1$ (Eichler & Waxman,
2005), where $m_{e}$ ($m_{p}$) is the electron (proton) mass. The external-
forward shock flux depends on $\xi$ and $\epsilon_{B}$ as
$f_{\nu}\propto\xi^{2-p}\epsilon_{B}^{(p+1)/4}$ (Leventis et al., 2012). From
this dependence, we find that the $\epsilon_{B}$ upper limit
$\propto\xi^{4(p-2)/(p+1)}$. Thus, including $\xi$ will decrease the values of
our $\epsilon_{B}$ upper limits. Values of $\xi$ have not been determined for
GRB external-forward shocks so we cannot quantify by how much the
$\epsilon_{B}$ upper limit values will decrease. Including $f$ will also
decrease the values of the $\epsilon_{B}$ upper limits since
$\epsilon_{B}^{\prime}=f\epsilon_{B}$ and $f<1$. Like $\xi$, values of $f$
have also not been determined from afterglow observations. The largest effect
$f$ can have on the $\epsilon_{B}$ upper limit values is decrease them by a
factor of $m_{p}/m_{e}\sim 2000$. For the remainer of this paper, we will be
conservative and continue to assume $\xi=f=1$, but we should keep in mind that
considering $\xi$ and $f$ will decrease the values of the $\epsilon_{B}$ upper
limits.
## 4\. Measurement Of $\epsilon_{B}$ With Optical Light Curves
### 4.1. $\epsilon_{B}$ Determination With Optical Data
The light curves we consider for our optical sample decline with a temporal
decay index $\alpha\sim 1$ from early times, $\sim 10^{2}-10^{3}\mbox{ sec}$,
as expected for the external-forward shock emission (see Section 4.2). Since
the light curves of these bursts are likely dominated by the external-forward
shock, this means that the observed optical flux is an actual measurement of
the external-forward shock flux, that is101010At late times, $\sim
10^{3}-10^{5}$ sec, many X-ray light curves decline with $\alpha\sim 1$: the
“normal” decline. If this segment arises from the external-forward shock, then
$\epsilon_{B}$ can be determined as described in this Section for our optical
sample. However, this is not straightforward since energy injection (as
evidenced by the plateau phase) should be considered. Also, one can attempt to
use the upper limit on $\epsilon_{B}$, found in Section 3, to calculate a
lower limit on $\nu_{c}$ during the normal decline and compare this to the
observed spectral regime ($\mbox{1 keV}<\nu_{c}$ or $\mbox{1 keV}>\nu_{c}$).
However, there are two difficulties with this $\nu_{c}$ consistency check: 1.
several “closure relations” are simultaneously satisfied within 2-$\sigma$ for
most of our sample (due to large errors in the spectral and temporal indices
at late times, which can be found in Evans et al. (2007, 2009)). 2\. for the
bursts that uniquely satisfy one closure relation, $\nu_{c}$ cannot be
determined precisely since Klein-Nishina suppression is weaker at late times
and $\nu_{c}$ depends strongly on $n$ when Compton-$Y$ is considered.
$f_{\mathrm{obs}}=f_{\mathrm{ES}}(E,n,\epsilon_{e},\epsilon_{B},p,s).$ (11)
Later in this section we will use this condition to determine $\epsilon_{B}$
for the bursts in our optical sample. We want to stress that we determine
$\epsilon_{B}$ for the optical sample. This is in contrast to the X-ray
sample, which only allowed us to determine an upper limit on $\epsilon_{B}$.
### 4.2. The Optical Sample
Our optical sample consists of 35 GRBs. $33/35$ of the bursts triggered Swift
and the remaining two bursts, 050502A and 080603A, were detected by INTEGRAL
(INTErnational Gamma-Ray Astrophysics Laboratory, Winkler et al. 2003). Table
LABEL:optical_properties_table shows properties of our optical sample. With a
few exceptions, most of the GRBs in our optical sample have a known
redshift111111The redshifts were taken from NASA’s Swift GRB Table and Lookup
website (exceptions are GRB 071003 (Perley et al., 2008) and GRBs 050502A and
071025 (J. Greiner’s website)). Three GRBs (050721, 070420, 060111B) in our
optical sample do not have a known redshift. For GRB 050721 and GRB 070420, we
assumed the average redshift of 2.5 for Swift GRBs (Gehrels et al., 2009). The
redshift for GRB 060111B was approximated as 2 in Stratta et al. (2009)..
Table 4Optical Sample Properties GRB | $z$ | $d_{L28}$ | Fluence | $E^{\gamma}_{\mathrm{iso,}52}$ | $\alpha_{O}$ | Ref. | $t_{2}$ | $f_{2\mathrm{eV}}$ | $\mbox{log}_{10}(\epsilon_{B})$
---|---|---|---|---|---|---|---|---|---
| | | $[\times 10^{-6}\mbox{ ergs}/\mbox{cm}^{2}]$ | | | | | [mJy] | ($s=0$)
050401 | 2.9 | 7.67 | 8.22 | 15.55 | $0.80\pm 0.03$ | [1] | 0.72 | 0.3 | -5.5
050502A | 3.793 | 10.59 | 1.4 | 4.12 | $1.16\pm 0.03$ | [2] | 1 | 5 | -4.5
050525A | 0.606 | 1.10 | 15.3 | 1.45 | $1.12\pm 0.05$ | [3] | 34.56 | 0.5 | -4.3
050721 | 2.5 | 6.40 | 3.62 | 5.32 | $1.29\pm 0.06$ | [4] | 20 | 0.2 | -5.0
050730 | 3.97 | 11.19 | 2.38 | 7.53 | $0.89\pm 0.05$ | [17] | 7.5 | 0.57 | -3.9
050802 | 1.71 | 4.00 | 2.00 | 1.49 | $0.82\pm 0.03$ | [1] | 3.6 | 0.5 | -3.5
051111 | 1.55 | 3.54 | 4.08 | 2.52 | $1.00\pm 0.02$ | [2] | 30 | 0.4 | -3.8
051221A | 0.5465 | 0.97 | 1.15 | 0.09 | $0.96\pm 0.03$ | [5] | 100 | 0.02 | -3.4
060111B | 2 | 4.86 | 1.60 | 1.58 | $1.18\pm 0.05$ | [6] | 2 | 0.4 | -5.2
060210 | 3.91 | 10.99 | 7.66 | 23.65 | $1.03\pm 0.06$ | [2] | 10 | 0.1 | -6.0
060418 | 1.49 | 3.37 | 8.33 | 4.78 | $1.13\pm 0.02$ | [7] | 2 | 8 | -4.6
060607A | 3.082 | 8.25 | 2.55 | 5.34 | $1.20\pm 0.03$ | [7] | 2 | 10 | -4.2
060904B | 0.703 | 1.32 | 1.62 | 0.21 | $1.00\pm 0.18$ | [17] | 5.5 | 0.58 | -3.5
060908 | 2.43 | 6.18 | 2.80 | 3.91 | $1.05^{+0.03}_{-0.03}$ | [8] | 2 | 2 | -4.5
060927 | 5.6 | 16.81 | 1.13 | 6.08 | $1.21\pm 0.06$ | [2] | 0.5 | 2 | -5.5
061007 | 1.26 | 2.74 | 44.4 | 18.48 | $1.70\pm 0.02$ | [7] | 2 | 50 | -6.0
061110B | 3.44 | 9.42 | 1.33 | 3.34 | $1.64\pm 0.08$ | [2] | 20 | 0.02 | -5.9
061121 | 1.314 | 2.88 | 13.7 | 6.19 | $0.82\pm 0.02$ | [7] | 4 | 0.5 | -4.7
061126 | 1.159 | 2.47 | 6.77 | 2.39 | $0.89\pm 0.02$ | [2] | 10 | 0.2 | -4.6
070318 | 0.84 | 1.65 | 2.48 | 0.46 | $0.96\pm 0.03$ | [7] | 20 | 0.2 | -3.7
070411 | 2.954 | 7.84 | 2.70 | 5.27 | $0.92\pm 0.04$ | [2] | 20 | 0.07 | -4.7
070420 | 2.5 | 6.40 | 14.0 | 20.56 | $0.81\pm 0.04$ | [2] | 3 | 0.8 | -4.7
070714B | 0.92 | 1.85 | 0.72 | 0.16 | $0.83\pm 0.04$ | [2] | 10 | 0.03 | -3.7
071003 | 1.6 | 3.69 | 8.3 | 5.45 | $1.466\pm 0.006$ | [9] | 0.6 | 20 | -5.7
071025 | 5.2 | 15.41 | 6.5 | 31.26 | $1.27\pm 0.04$ | [10] | 20 | 0.02 | -6.8
071031 | 2.692 | 7.00 | 0.9 | 1.5 | $0.97\pm 0.06$ | [11] | 10.5 | 0.4 | -3.4
071112C | 0.823 | 1.61 | 3.00 | 0.53 | $0.95\pm 0.02$ | [12] | 10.5 | 0.003 | -6.3
080603A | 1.688 | 3.94 | 1.1 | 0.80 | $0.99\pm 0.07$ | [13] | 30 | 0.1 | -3.6
080607 | 3.036 | 8.10 | 24.0 | 49.04 | $1.65$ | [14] | 3 | 0.2 | -8.0
080721 | 2.591 | 6.68 | 12.0 | 18.74 | $1.22\pm 0.01$ | [5] | 3 | 10 | -5.0
080810 | 3.35 | 9.13 | 4.60 | 11.06 | $1.23\pm 0.01$ | [7] | 3 | 30 | -3.9
080913 | 6.7 | 20.72 | 0.56 | 3.92 | $1.03\pm 0.02$ | [15] | 10 | 0.02 | -5.2
081008 | 1.967 | 4.76 | 4.30 | 4.13 | $0.96\pm 0.03$ | [16] | 2 | 3 | -4.1
090313 | 3.375 | 9.21 | 1.40 | 3.41 | $1.25\pm 0.08$ | [17] | 20 | 2 | -3.4
090418A | 1.608 | 3.71 | 4.60 | 3.05 | $1.21\pm 0.04$ | [7] | 2 | 0.8 | -5.5
Table 4Optical Sample Properties (Continued)
With the exception of the only short GRB in our optical sample, GRB 051221A,
all the optical light curves in our sample decline before 3500 seconds.
Considering early times has the advantage of minimizing possible energy
injection. Our optical sample can be separated into 4 different subgroups,
depending on the temporal behavior of the light curve before the $\alpha\sim
1$ decay as follows.
1. 1.
Light Curves With A Peak At Early Times: The light curves of this subgroup are
characterized by a power law rise, reaching a peak, and then a power law
decline with $\alpha_{O}\sim 1$ ($\sim$43% of our sample). The peak of the
light curve is believed to be due to the deceleration time. For the bursts in
this subgroup, we show the temporal decay of the optical light curve after the
peak and the flux and the time of the second data point after the peak in
Table LABEL:optical_properties_table. We take the second data point to be
confident that the optical light curve is declining.
2. 2.
Single Power Law Decay From Early Times: In this subgroup, the optical light
curve shows a decline as a single power law with $\alpha\sim 1$ from the
beginning of the observations ($\sim$ 40% of our sample). We display the
temporal decay of the optical light curve and the time and the flux of the
second data point observed in Table LABEL:optical_properties_table.
3. 3.
Optical Light Curves With Plateaus At Early Times: The optical light curves of
3 bursts in our optical sample (GRBs 050525A, 060210, and 070411) display
plateaus at early times. The plateaus in our optical sample are short, with
the longest plateau lasting under 3500 seconds. After the plateau ends, the
light curves of these 3 bursts show a decay $\alpha_{O}\sim 1$, as expected
for the external forward-shock emission. In Table
LABEL:optical_properties_table, for these 3 bursts, we show the temporal decay
after the plateau and the time and the flux of the second data point after the
plateau.
4. 4.
Light Curves With Possible Reverse Shock Emission At Early Times: 3 GRBs in
our optical sample (060111B, 060908, and 061126) show possible emission from
the reverse shock. The light curves in this subgroup show an initial steep
decline at early times, characteristic of the reverse shock, and then
transition to a more shallow decay of $\alpha\sim 1$ that is more typical for
the external-forward shock emission. For these GRBs, in Table
LABEL:optical_properties_table we show the temporal decay of the light curve
and the time and flux of the second data point after the possible reverse
shock emission.
### 4.3. Optical External-Forward Shock Spectral Regime and Afterglow
Parameter Assumptions
When referring to the optical band, we will use 2 eV since most of the light
curves in our optical sample are either plotted at 2eV or were observed in the
R filter. As we did with the X-ray sample, we will only consider the slow
cooling ordering of the synchrotron characteristic frequencies,
$\nu_{a}<\nu_{i}<\nu_{c}$. Because the optical light curve is declining at the
time we are considering, the optical band must be above $\nu_{i}$ at this
time. In Section 3.4, we argued that the X-ray band is between $\nu_{i}$ and
$\nu_{c}$ at the end of the steep decline at a few 100 sec; therefore, the
optical band must also be in this spectral regime at the early times $(\sim
10^{2}-10^{3})$ sec we are considering. The expression we will use to
determine the optical external-forward shock flux is also Equation 7; however,
we will have an equality (instead of an inequality), we replace
$f_{\mathrm{1keV,EoSD}}$ with $f_{\mathrm{2eV}}$ (which represents the
specific flux observed at 2 eV), and use $\nu_{14}$ corresponding to 2 eV.
The other afterglow parameters are determined as in Section 3.2:
$\epsilon_{e}=0.2$ and with $z$ and the fluence121212For the two bursts
detected by INTEGRAL, the fluence in Table LABEL:optical_properties_table is
in the 20-200 keV band of the instrument IBIS (Imager on-Board the INTEGRAL
Satellite, Ubertini et al. (2003)): GRB 050502A (Gotz & Mereghetti, 2005) and
GRB 080603A (Guidorzi et al., 2011)., we obtain $E^{\gamma}_{\mathrm{iso}}$
and use $E=5E^{\gamma}_{\mathrm{iso}}$. As with our X-ray sample, we will
display our $\epsilon_{B}$ results with $p=2.2$, 2.4, and 2.8. We can also
determine $p$ by using the temporal decay of the optical light curve,
$\alpha_{O}$, which is shown in Table LABEL:optical_properties_table for each
burst (optical spectrum is not always available, so we cannot use the closure
relations for the optical sample). In order to have $p>2$ for all of the
bursts in our optical sample, we only consider a constant density medium when
determining $p$ with $\alpha_{O}$ ($\alpha_{O}=3(p-1)/4$ for $s=0$ and
$\alpha_{O}=(3p-1)/4$ for $s=2$). Lastly, as we did for our X-ray sample, to
keep the density-dependence, we will plot the quantity
$\epsilon_{B}n_{0}^{2/(p+1)}$ ($\epsilon_{B}A_{*-1}^{4/(p+1)}$) for $s=0$
($s=2$).
### 4.4. $\epsilon_{B}$ Results For Optical Sample
$\begin{array}[]{c}\begin{array}[]{cc}\includegraphics[scale={0.6}]{optical_p_2_2_histogram.eps}&\includegraphics[scale={0.6}]{optical_p_2_4_histogram.eps}\end{array}\\\
\begin{array}[]{cc}\includegraphics[scale={0.6}]{optical_p_2_8_histogram.eps}&\includegraphics[scale={0.6}]{s_0_optical_histogram_with_temporal_p.eps}\end{array}\end{array}$
Figure 3.— The Top-Left, Top-Right, and Bottom-Left panels show the histograms
of measurements assuming $p=2.2$, $p=2.4$, and $p=2.8$ respectively for all of
the GRBs in our optical sample. The filled-in (un-filled) histograms show
measurements of the quantity $\epsilon_{B}n_{0}^{2/(p+1)}$
($\epsilon_{B}A_{*-1}^{4/(p+1)}$) assuming all the GRBs in our optical sample
are described by a constant density (wind) medium. Bottom-Right Panel: The
filled-in histogram shows the measurements on the quantity
$\epsilon_{B}n_{0}^{2/(p+1)}$ with $p$ determined from $\alpha_{O}$. The un-
filled histogram shows measurements on the quantity
$\epsilon_{B}n_{0}^{2/(p+1)}$, assuming $p=2.4$ for all of the bursts in our
optical sample (this histogram was also shown in the Top-Right panel).
We display the results for the measurements (from Equation 7) on the quantity
$\epsilon_{B}n_{0}^{2/(p+1)}$ ($\epsilon_{B}A_{*-1}^{4/(p+1)}$) for $s=0$
($s=2$) assuming all the GRBs in our optical sample have $p=2.2$, 2.4, and 2.8
in the Top-Left, Top-Right, and Bottom-Left panels of Figure 3, respectively.
Two histograms are shown in each panel, one for $s=0$ and the other for $s=2$.
We also use $\alpha_{O}$ to determine $p$ (assuming $s=0$) and compare the
results to the ones obtained with $p=2.4$ and $s=0$ (Bottom-Right panel of
Figure 3). In Table 3, we display a summary of the mean and median values of
the measurements of the quantity $\epsilon_{B}n_{0}^{2/(p+1)}$
($\epsilon_{B}A_{*-1}^{4/(p+1)}$) for $s=0$ ($s=2$) for each histogram. For
the remainder of this section, we assume a standard $n_{0}=1$ ($A_{*-1}=1$)
for $s=0$ ($s=2$) when discussing our results for the $\epsilon_{B}$
measurements for our optical sample.
For a constant density (wind) medium, the mean and median $\epsilon_{B}$
measurements are $\sim\mbox{ few}\times 10^{-5}$ ($\sim\mbox{ few}\times
10^{-6}$) The mean and median $\epsilon_{B}$ measurements only change by a
factor of a few when assuming a different value of $p$. To determine if
assuming a standard $p=2.4$, as opposed to determining $p$ for each burst from
$\alpha_{O}$, significantly affects the distribution of $\epsilon_{B}$
measurements, we compared the two histograms in the Bottom-Right panel of
Figure 3 with a Kolmogorov-Smirnov (KS) test. The null hypothesis of the KS
test is that the two histograms are drawn from the same distribution. We test
this null hypothesis at the 5% significance level. The KS test confirmed the
null hypothesis that the two histograms are consistent with being drawn from
the same distribution.
As with the $\epsilon_{B}$ upper limits from X-ray data, the mean and median
$\epsilon_{B}$ measurements decrease by about an order of magnitude when
assuming a wind medium as opposed to a constant density medium. Compared to
the distribution of $\epsilon_{B}$ upper limits we attained from X-ray data,
the $\epsilon_{B}$ measurements from optical data show a much wider
distribution. For a constant density (wind) medium, the $\epsilon_{B}$
measurements range from $\epsilon_{B}\sim 10^{-8}-10^{-3}$ ($\epsilon_{B}\sim
10^{-9}-10^{-3}$). Also, since we used the same equation (Equation 7) to find
both the upper limits on $\epsilon_{B}$ with X-ray data and the $\epsilon_{B}$
measurements with optical data, the discussion at the end of Section 3.5 on
how the uncertainty in the afterglow parameters and the parameters $\xi$ and
$f$ can affect the distribution of $\epsilon_{B}$ upper limits also applies to
the distributions of $\epsilon_{B}$ measurements we presented in this section.
In addition, since $\xi$ and $f$ are less than unity, including these two
parameters will mean that our $\epsilon_{B}$ measurements are effectively
upper limits on $\epsilon_{B}$.
### 4.5. Comparison Of Our Results On $\epsilon_{B}$ To Previous Studies
We performed a KS test between our optical $\epsilon_{B}$ measurements and the
results from previous studies on $\epsilon_{B}$ (Figure 1). For our
$\epsilon_{B}$ results, we used the optical $\epsilon_{B}$ measurements with
$n=1\mbox{ cm}^{-3}$ and $p$ determined from $\alpha_{O}$ (filled-in histogram
in Bottom-Right panel of Figure 3; the $\epsilon_{B}$ values are shown in
Table LABEL:optical_properties_table). The result of the KS test is that the
null hypothesis is rejected. The P-value, which measures the probability that
the null hypothesis is still true, is $2.1\times 10^{-9}$. This result shows
that the rejection of the null hypothesis is statistically significant. It is
not surprising that the null hypothesis was rejected. The distribution from
the previous studies is very inhomogeneous, with the values for $\epsilon_{B}$
being drawn from many different studies with different methodologies. Also,
comparing the histogram in Figure 1 to the filled-in histogram in the Bottom-
Right panel of Figure 3, a couple of significant differences can be seen. The
range for the histogram of $\epsilon_{B}$ values found in the literature is
$\sim 10^{-5}-10^{-1}$, whereas the range for our $\epsilon_{B}$ results is
$\sim 10^{-8}-10^{-3}$. The mean and median values for these two histograms
are also significantly different. The mean (median) value for the
$\epsilon_{B}$ histogram from the literature, $6.3\times 10^{-2}$ ($1.4\times
10^{-2}$), is a factor $\sim 700$ ($\sim 600$) times larger than the mean
(median) $\epsilon_{B}$ value of the histogram with our results, which is
$9.5\times 10^{-5}$ ($2.4\times 10^{-5}$).
One assumption that is commonly made in afterglow modeling studies is
equipartition between $\epsilon_{e}$ and $\epsilon_{B}$. As we discussed in
Section 2, the results for $\epsilon_{e}$ from the literature and the results
from recent simulations of relativistic collisionless shocks support
$\epsilon_{e}\sim 0.2$. From this result, many works assume $\epsilon_{B}\sim
10^{-2}-10^{-1}$. However, there is no physical argument to expect
equipartition. Our distribution of $\epsilon_{B}$ upper limits and
measurements, although wide, supports that there is no equipartition between
electron and magnetic energies because none of the $\epsilon_{B}$ upper limits
or measurements in our samples has a value larger than $\epsilon_{B}\sim\mbox{
few }\times 10^{-3}$. Another source of error that can lead to differences in
$\epsilon_{B}$ values is differences in the determination of the spectral
regime for the optical band. We took it to be between $\nu_{i}$ and $\nu_{c}$,
but it is also possible for the optical band to be above $\nu_{c}$ at late
times (e.g. Panaitescu & Kumar, 2002; Cenko et al., 2010). Another source for
error is energy injection. We did not consider energy injection as a source of
error because only $3/35$ of the bursts in our optical sample show plateaus
(and these plateaus are short). The X-ray and optical light curves of many
bursts show plateaus and in these cases energy injection needs to be
considered. Also, errors in our determination of fluxes and times from X-ray
and optical light curves can also lead to small errors in $\epsilon_{B}$. In
summary, the main assumption we made when determining $\epsilon_{B}$ is an
efficiency of $\sim 20\%$ in the conversion of kinetic energy to gamma-ray
energy, and we did not assume equipartition between $\epsilon_{e}$ and
$\epsilon_{B}$. Different authors have made different assumptions that can
have a large effect on the results for $\epsilon_{B}$.
Lastly, for a few bursts, we checked if our method of determining
$\epsilon_{B}$ is consistent with the values determined for $\epsilon_{B}$
with other techniques. GRBs 980519 and 990123, discussed in the afterglow
modelling study of Panaitescu & Kumar 2002, have optical light curves that
decline as a power-law before the jet-break. The optical band for both of
these bursts was determined to be in the spectral regime $\nu_{i}<\mbox{ 2
eV}<\nu_{c}$. Applying our technique to find a $\epsilon_{B}$ measurement for
both of these bursts and using the value of $n$ reported in Panaitescu & Kumar
2002 for both of these bursts, we find that these bursts have
$\epsilon_{B}\sim 10^{-5}$, consistent with the results reported in Panaitescu
& Kumar 2002 for both of these bursts within a factor of a few. The small
differences in $\epsilon_{B}$ values can be accounted for by differences in
the coefficients used for the external-forward shock flux.
## 5\. GRBs In Common To Both Our X-ray And Optical Samples
14 bursts we considered are both in our X-ray and in our optical sample (GRB
number is in bold in the first columns of Tables LABEL:x-ray_properties_table
and LABEL:optical_properties_table). In this section, we verify for these
bursts that the X-ray $\epsilon_{B}$ upper limit is above the optical
$\epsilon_{B}$ measurement. For the optical data, we will use the
$\epsilon_{B}$ measurements with $p$ determined from $\alpha_{O}$ and
$n=1\mbox{ cm}^{-3}$ (filled-in histogram in the Bottom-Right panel of Figure
3). For the X-ray data, in this section, we will also consider $n=1\mbox{
cm}^{-3}$ ($s=0$) and we will use the value of $p$ determined from
$\alpha_{O}$ 13131310 out of 14 of these bursts have optical data before 1000
seconds. For these 10 bursts, we can use the optical data to check if they
satisfy the assumption we made in Section 3.1,
$t_{\mathrm{dec}}<t_{\mathrm{EoSD}}$. 9 of these 10 GRBs do satisfy this
assumption; for the remaining GRB (080721), we are not able to check this
assumption because the first optical observation (at 100 sec.) is after the
end of the steep decline ($t_{\mathrm{EoSD}}=70$ sec).. The comparison between
the X-ray $\epsilon_{B}$ upper limits and the optical $\epsilon_{B}$
measurements is shown in Figure 4. This plot shows that all the X-ray
$\epsilon_{B}$ upper limits are above the optical $\epsilon_{B}$ measurements.
Figure 4.— Comparison of the $\epsilon_{B}$ upper limits from X-ray data to
the $\epsilon_{B}$ measurements from optical data. The 14 dots correspond to
the 14 GRBs that are both in our X-ray and optical samples. The straight line
indicates where the $\epsilon_{B}$ measurements are equal to the
$\epsilon_{B}$ upper limits.
## 6\. $E$ vs. $\epsilon_{B}$ Correlation?
We now use the $\epsilon_{B}$ measurements from our optical sample to
determine if there is a correlation between $E$ and $\epsilon_{B}$. The 3
bursts studied in Kumar & Barniol Duran 2009, 2010 have small values of
$\epsilon_{B}$, consistent with shock compression of a seed magnetic field
$B_{0}\sim\mbox{ few}\times 10\mu\mbox{G}$. One property that distinguishes
these 3 GRBs is that they were particularly energetic, with
$E_{\mathrm{iso}}^{\gamma}\sim 10^{53}-10^{54}\mbox{ ergs}$. Could the large
energy intrinsic to these 3 bursts explain why these 3 bursts have low
$\epsilon_{B}$ values? We investigate this possibility in Figure 5 by plotting
the values of $E$ and $\epsilon_{B}$ for all the bursts in our optical sample.
Figure 5.— We plot the values of $E$ and the measurements of $\epsilon_{B}$ to
determine if they are correlated. The 35 points represent the GRBs in our
optical sample and the straight line is the best fit line:
$\mbox{log}_{10}(\epsilon_{B})=-1.02\mbox{log}_{10}(E_{53})-4.51$, with the
slope of the line being $-1.02\pm 0.23$ and the y-intercept of the line being
$-4.51\pm 0.16$. The correlation coefficient of the fit is 0.62 and the
P-value of the correlation is $1.2\times 10^{-4}$ (3.8$\sigma$ significance).
The $\epsilon_{B}$ measurements are for $n=1\mbox{ cm}^{-3}$ and $p$
determined from $\alpha_{O}$ (shown in the filled-in histogram in the Bottom-
Right panel of Figure 3 and Table LABEL:optical_properties_table) and the
values of $E$ were determined by assuming an efficiency $\sim 20\%$ for all
the GRBs in our optical samples.
For the energy of each burst, we assumed an efficiency of $\sim 20\%$
($E=5E^{\gamma}_{\mathrm{iso}}$) in the conversion of the kinetic energy of
the jet to gamma-ray radiation; the observables involved in calculating $E$
are the gamma-ray fluence and $z$. For the $\epsilon_{B}$ measurements, we
used the values with $p$ determined from $\alpha_{O}$ and $n=1\mbox{
cm}^{-3}$. The observables involved in determining $\epsilon_{B}$ are the
observed specific optical flux and the time. In Figure 5, we also show the
best fit line141414From Equation 7, it may be expected that $E$ and
$\epsilon_{B}$ are correlated. This is not necessarily true because each burst
has a different value for the observed quantities $f_{\nu}$ and $t$. In
addition, since $p$ was determined from $\alpha_{O}$, each burst has a
different $p$. We also checked if $E$ and $\epsilon_{B}$ were correlated when
assuming a fixed $p=2.4$. From Equation 9, we would expect the slope of the
line to be $-1.59$. The best fitting line for $p=2.4$ is
$\mbox{log}_{10}(\epsilon_{B})=-0.77\mbox{log}_{10}(E_{53})-4.53$, with the
slope of the line being $-0.77\pm 0.24$ and the y-intercept being $-4.53\pm
0.16$. The slope of this line is more than $3\sigma$ away from $-1.59$,
showing that there is no expected correlation between $E$ and $\epsilon_{B}$
for our methodology of determining $\epsilon_{B}$.. This fit shows that an
increase in $E$ leads to a decrease in $\epsilon_{B}$. The correlation
coefficient of the fit is $0.62$, indicating that there is a weak correlation
between $E$ and $\epsilon_{B}$. A value of the correlation coefficient close
to 1 would indicate a strong correlation. The P-value of the correlation is
$1.2\times 10^{-4}$ (3.8$\sigma$ significance), indicating that there is a
small probability that the correlation occurred by chance.
Although many points show large deviations from the best fit line, the scatter
of the points may be reduced or increased by the uncertainty in the afterglow
parameters. An error in the efficiency would affect the values of $E$ and an
error in $\epsilon_{e}$, the efficiency, or $n$ would affect the
$\epsilon_{B}$ measurements (see Section 3.5 for a discussion on how the
$\epsilon_{B}$ measurements would be affected when an error in a afterglow
parameter is made). It is possible that the uncertainty in $E$ and
$\epsilon_{B}$ can reduce or increase the scatter and make the correlation
between $E$ and $\epsilon_{B}$ stronger or weaker.
## 7\. Magnetic Field Amplification Factor For X-ray And Optical Results
In Sections 3-5, we presented our results for the strength of the magnetic
field downstream of the shock front in terms of the afterglow parameter
$\epsilon_{B}$. If shock compression was the only mechanism amplifying the
ambient magnetic field (assuming a standard $B_{0}\sim\mbox{few }\mu\mbox{G}$
and a standard $n=1\mbox{ cm}^{-3}$), then $\epsilon_{B}\sim 10^{-9}$ is
expected. Most of the bursts in our distributions of $\epsilon_{B}$ upper
limits and measurements have values larger than $\epsilon_{B}\sim 10^{-9}$.
These results suggest that amplification of the magnetic field, in addition to
shock compression, is needed to explain the afterglow observations. In this
section, we will present our results in terms of an amplification factor,
which quantifies the amplification that is needed, beyond shock compression,
to explain the observations.
If shock compression were the only mechanism amplifying the seed magnetic
field $B_{0}$, then $B=4\Gamma B_{0}$. To quantify how much additional
amplification of the ambient magnetic field is needed, beyond shock
compression, we define the amplification factor, $AF$, as
$AF\equiv\frac{B}{4B_{0}\Gamma}.$ (12)
$AF$ is a constant that satisfies $AF\geq 1$ since $B\geq 4B_{0}\Gamma$.
$AF=1$ means than the observed $B$ is consistent with the only amplification
arising from seed magnetic field shock compression. The expression for
$\epsilon_{B}$ is $\epsilon_{B}=B^{2}/32\pi m_{p}c^{2}n\Gamma^{2}$. With the
definition for $AF$, $\epsilon_{B}$ is
$\epsilon_{B}=(AF)^{2}\times\frac{B_{0}^{2}}{2\pi nm_{p}c^{2}}.$ (13)
We note that $\epsilon_{B}$ is given by $(AF)^{2}$ times the $\epsilon_{B}$ we
would get if shock compression were the only mechanism amplifying the magnetic
field.
We will now use Equation 13 and our previous results for the X-ray
$\epsilon_{B}$ upper limits to determine an upper limit on $AF$. In Section
3.5, if we assumed a standard $n=1\mbox{ cm}^{-3}$, we were able to attain an
upper limit on $\epsilon_{B}$ for a constant density medium 151515We will only
consider a constant density medium when displaying the results for the
amplification factor. We will show in Equation 15 that $AF$ has a weak
dependence on the density.. We will refer to these $\epsilon_{B}$ upper limits
as $\overline{\epsilon_{B}}(p,n_{0}=1)$. In the notation
$\overline{\epsilon_{B}}(p,n_{0}=1)$, the bar over $\epsilon_{B}$ signifies
that this is an upper limit on $\epsilon_{B}$, the $p$ in the parenthesis
shows that the $\epsilon_{B}$ upper limit depends on the value of $p$ we used,
and the $n_{0}=1$ shows that we assumed $n=1\mbox{ cm}^{-3}$. With this
notation, we can keep the dependence of the $\epsilon_{B}$ upper limit on $n$
(see Equation 7):
$\epsilon_{B}<\frac{\overline{\epsilon_{B}}(p,n_{0}=1)}{n_{0}^{2/(p+1)}}.$
(14)
Combining Equation 13 and Equation 14, the expression for the upper limit on
$AF$ is
$AF<\frac{1.0\times
10^{4}[\overline{\epsilon_{B}}(p,n_{0}=1)]^{1/2}}{B_{0,10\mu\mathrm{G}}n_{0}^{(1-p)/(2p+2)}},$
(15)
where $B_{0,10\mu\mathrm{G}}\equiv B_{0}/10\mu\mbox{G}$. One advantage to
expressing the results of the magnetic field downstream of the shock front in
terms of $AF$ is that $AF$ depends weakly on $n$. For a standard $p=2.4$,
$AF\propto n^{0.21}$. On the other hand, there is a strong dependence on
$B_{0}$, $AF\propto B_{0}^{-1}$.
### 7.1. Amplification Factor Upper Limit For Our X-ray Sample
We will now show the results for the $AF$ upper limits we obtained from
Equation 15. Since the amplification factor has a weak dependence on the
density, we will assume a standard $n_{0}=1$ when displaying the results for
the $AF$ upper limits. When plotting the results for $AF$, we will keep the
dependence on $B_{0}$ and plot the quantity $(AF)B_{0,10\mu\mathrm{G}}$. In
the left and right panels of Figure 6, we show the upper limits on the
quantity $(AF)B_{0,10\mu\mathrm{G}}$ for a fixed $p=2.4$ and $p=2.2$, 2.8,
where the values of $\overline{\epsilon_{B}}(p,n_{0}=1)$ used in Equation 15
were shown in Top-Right, Top-Left, and Bottom panels of Figure 2,
respectively. For the remainder of this section, we will assume
$B_{0}=10\mu\mbox{G}$ when discussing the results for the $AF$ upper limits.
The mean and median values of the $AF$ upper limits are summarized in Table 5.
The mean (median) $AF$ upper limits range from $AF\sim 60-80$ ($AF\sim
40-60$). The $AF$ upper limit histograms show a wide distribution, with a
range of $\sim 10$ to $\sim 300$. To determine if assuming a different value
of $p$ has a significant effect on the distribution of $AF$ upper limits, we
performed a KS test between the histograms in the right panel of Figure 6. The
KS test confirmed the null hypothesis, leading us to conclude that the $AF$
upper limit results are not sensitive to the value of $p$ we assume.
X-ray | | | |
---|---|---|---|---
$(s=0)$ | $p=2.2$ | $p=2.4$ | $p=2.8$ |
Mean | $84$ | $67$ | $62$ |
Median | $63$ | $51$ | $44$ |
Opt. | | | |
$(s=0)$ | $p=2.2$ | $p=2.4$ | $p=2.8$ | $p$ from $\alpha_{\mathrm{O}}$
Mean | $130$ | $71$ | $36$ | $71$
Median | $100$ | $56$ | $23$ | $48$
Table 5Mean and median $AF$ values for the X-ray (upper limits on $AF$) and
optical (measurements of $AF$) histograms shown in Figures 6 and 7. All the
labels are the same as in Table 3. A constant density medium with $n=1\mbox{
cm}^{-3}$ (the amplification factor has a weak dependence on the density, see
Section 7 and Equation 15) and a seed magnetic field $B_{0}=10\mu\mbox{G}$
were assumed for all the bursts in our X-ray and optical samples.
$\begin{array}[]{c}\begin{array}[]{cc}\includegraphics[scale={0.6}]{amplification_factor_x_ray_p_2_4.eps}&\includegraphics[scale={0.6}]{amplification_factor_x_ray_p_2_2_p_2_8.eps}\end{array}\end{array}$
Figure 6.— Left: Upper limits on the quantity $(AF)B_{0,10\mu\mathrm{G}}$ for
our X-ray sample assuming $p=2.4$. Right: Upper limits on the quantity
$(AF)B_{0,10\mu\mathrm{G}}$ for our X-ray sample assuming $p=2.2$ and $p=2.8$.
A fixed $n=1\mbox{ cm}^{-3}$ was assumed for all of the histograms (the
precise value of $n$ is unimportant since $AF$ has a weak dependence on $n$,
see Section 7 and Equation 15).
We now discuss how an error in each of the afterglow parameters can affect our
results for the $AF$ upper limits. For this discussion, we will assume
$p=2.4$. From Equation 15,
$\mbox{AF}\propto\sqrt{\overline{\epsilon_{B}}(p=2.4,n_{0}=1)}B_{0}^{-1}n^{0.2}$
and from Equation 9, $\epsilon_{B}\propto\epsilon_{e}^{-1.6}E^{-1.6}n^{-0.6}$.
From these two expressions,
$\mbox{AF}\propto\epsilon_{e}^{-0.8}E^{-0.8}n^{0.2}B_{0}^{-1}$. We note that
compared to the $\epsilon_{B}$ upper limit (Equation 9), the $AF$ upper limit
has a weaker dependence on $\epsilon_{e}$, $E$, and $n$. As we discussed at
the end of Section 3.5, a likely error in $\epsilon_{e}$ is a factor of $\sim
2$; this error in $\epsilon_{e}$ will translate into an error in the $AF$
upper limits by a factor of only $\sim 2$. For the energy, we assumed an
efficiency $\sim 20\%$ and a likely error in the efficiency is a factor $\sim
2-3$ (see Section 3.5); this error in the efficiency would lead to an error in
the $AF$ upper limits by only a factor $\sim 2$. One advantage to expressing
the results of the magnetic field downstream of the shock front in terms of
$AF$ is that $AF$ has a very weak dependence on $n$. An error in $n$ by a
factor $\sim 10^{3}$ (see Section 3.5) from our assumed $n=1\mbox{ cm}^{-3}$
will only lead to an error in the $AF$ upper limits by a factor $\sim 4$. The
price to pay for a weak $n$ dependence is a linear dependence on $B_{0}$, with
$\mbox{AF}\propto B_{0}^{-1}$. $B_{0}$ is an uncertain parameter that likely
varies from GRB environment to GRB environment and it is the largest source of
uncertainty for $AF$.
### 7.2. Amplification Factor Measurement For Our Optical Sample
As we discussed in Section 4.2, for our optical sample, we found a measurement
for $\epsilon_{B}$ instead of an upper limit. This will allow us to determine
a measurement for $AF$. To do this, we will use Equation 15, but in this case
we have an equality instead and we have $\epsilon_{B}(p,n_{0}=1)$ instead of
$\overline{\epsilon_{B}}(p,n_{0}=1)$. The notation $\epsilon_{B}(p,n_{0}=1)$
denotes the $\epsilon_{B}$ measurements for our optical sample from Section
4.4 if we assume a standard $n=1\mbox{ cm}^{-3}$. Also, as with the X-ray
sample, we only consider $s=0$ when calculating the $AF$ measurements and
assume a fixed $n=1\mbox{ cm}^{-3}$. In the left panel of Figure 7, we show
the results for the measurements on the quantity $(AF)B_{0,10\mu\mathrm{G}}$
for $p$ determined from $\alpha_{O}$ and also assuming $p=2.4$. In the right
panel of Figure 7, we show the measurements on the quantity
$(AF)B_{0,10\mu\mathrm{G}}$ for $p=2.2$ and $p=2.8$. For the remainder of this
section, we will assume $B_{0}=10\mu\mbox{G}$ when discussing the results for
the $AF$ measurements.
A summary of the mean and median $AF$ measurements for our optical sample is
shown in Table 5. To determine if assuming a standard $p=2.4$, as opposed to
determining $p$ from $\alpha_{O}$ for each burst, has a statistically
significant effect on the distribution of $AF$ measurements, we performed a KS
test between the two histograms in the left panel of Figure 7. The KS test
confirmed the null hypothesis. The mean (median) $AF$ measurements for the
optical histograms range from $\sim 40$ to $\sim 130$ ($\sim 20$ to $\sim
100$). Compared to the $AF$ upper limit histograms, the $AF$ measurement
histograms show a wider distribution, ranging from $AF\sim 1$ to $AF\sim
1000$. Also, since we used the same expression to determine the $AF$ upper
limits and measurements (Equation 15), the discussion at the end of Section
7.1 on how an error in one of the afterglow parameters can affect the $AF$
upper limits also applies to the $AF$ measurements161616$\xi$ and $f$ also
affect our $AF$ results. To account for $\xi$, since the $\epsilon_{B}$ upper
limit$/$measurement is $\propto\xi^{4(p-2)/(p+1)}$ and the $AF$ upper
limit$/$measurement is $\propto(\overline{\epsilon_{B}}(p,n_{0}=1))^{1/2}$,
$AF\propto\xi^{2(p-2)/(p+1)}$. To account for $f$, since
$AF\propto(\overline{\epsilon_{B}}(p,n_{0}=1))^{1/2}$, $AF\propto f^{1/2}$.
Thus, including $\xi$ and $f$ will decrease the values of the $AF$ upper
limits$/$measurements. Taking the lowest possible value for $f$, the $AF$
upper limits$/$measurements can decrease by up to a factor $\sim 40$. This
would make $\sim 50\%$($\sim 60\%$) of the bursts in our X-ray (optical)
sample consistent with shock compression. In addition, as with the
$\epsilon_{B}$ measurements, since $\xi$ and $f$ are less than unity,
including these two parameters will mean that our $AF$ measurements will
become upper limits on $AF$..
$\begin{array}[]{c}\begin{array}[]{cc}\includegraphics[scale={0.6}]{amplification_factor_optical_p_temporal_p_2_4.eps}&\includegraphics[scale={0.6}]{amplification_factor_optical_p_2_2_p_2_8.eps}\end{array}\end{array}$
Figure 7.— Left Panel: The filled-in (un-filled) histogram shows the
measurements on the quantity $(AF)B_{0,10\mu\mathrm{G}}$ for $p$ calculated
from $\alpha_{O}$ ($p=2.4$). Right Panel: The filled-in (un-filled) histogram
shows the measurements on the quantity $(AF)B_{0,10\mu\mathrm{G}}$ for $p=2.2$
($p=2.8$). A fixed $n=1\mbox{ cm}^{-3}$ was assumed for all of the histograms
(the precise value of $n$ is unimportant since $AF$ has a weak dependence on
$n$, see Section 7 and Equation 15).
## 8\. Discussion And Conclusions
In this work, we presented a systematic study on the magnetic field downstream
of the shock front for large samples of GRBs (60 in our X-ray sample and 35 in
our optical sample). We expressed the strength of the downstream magnetic
field in terms of both the afterglow parameter $\epsilon_{B}$ and a
amplification factor, denoted by $AF$, which quantifies how much
amplification, beyond shock compression of the seed magnetic field, is needed
to explain the downstream magnetic field. This is the first time a large and
systematic study has been carried out to study $\epsilon_{B}$ and to determine
how much amplification of the seed magnetic field is required by the
observations. For our X-ray (optical) sample, we determined an upper limit
(measurement) for both $\epsilon_{B}$ and $AF$. The upper limits on
$\epsilon_{B}$ and $AF$ for our X-ray sample were found from the constraint
that the observed flux at the end of the steep decline is greater than or
equal to the external-forward shock flux. This is a new method to constrain
$\epsilon_{B}$ that relies on the steep decline emission, which has been
observed by Swift for many GRBs. Our optical sample was restricted to light
curves that decline with $\alpha\sim 1$ from the early times $\sim
10^{2}-10^{3}$ sec, as expected for the external-forward shock emission. We
found the measurements for $\epsilon_{B}$ and $AF$ for our optical sample from
the condition that the observed flux is equal to the external-forward shock
flux.
The condition used for our X-ray (optical) sample was converted into an upper
limit (measurement) on the quantity $\epsilon_{B}n^{2/(p+1)}$ for $s=0$ or
$\epsilon_{B}A_{*-1}^{4/(p+1)}$ for $s=2$ by assuming a $\sim 20\%$ efficiency
in the conversion of kinetic energy to prompt gamma-ray radiation. To find an
upper limit (measurement) on $\epsilon_{B}$ for our X-ray (optical) sample, we
assumed a standard $n=1\mbox{ cm}^{-3}$ ($A_{*}=0.1$) for the density for a
constant density (wind) medium. A discussion on how the uncertainty in the
afterglow parameters affects our results for $\epsilon_{B}$ can be found at
the end of Section 3.5. The largest source of uncertainty for our results on
$\epsilon_{B}$ is the density, since the value of the density has been
observed to vary over many orders of magnitude and its precise value is not
known for each GRB (see Section 3.2). For the bursts that are both in our
X-ray and optical samples, we also applied a consistency check to make sure
our results for $\epsilon_{B}$ are correct (see Section 5).
From Table 3, for a constant density (wind) medium, most of the $\epsilon_{B}$
upper limit and measurement histograms have a median value $\sim\mbox{
few}\times 10^{-5}$ ($\sim\mbox{ few}\times 10^{-6}$). These results imply
that half of the bursts in both our X-ray and optical samples have a
$\epsilon_{B}$ value $\sim\mbox{ few}\times 10^{-5}$ or lower. Assuming
$n=1\mbox{ cm}^{-3}$ and $B_{0}\sim\mbox{few}\times\mu$G, shock compression is
only able to produce $\epsilon_{B}\sim 10^{-9}$. Although $\epsilon_{B}\sim
10^{-9}$ is 4 orders of magnitude lower than $\epsilon_{B}\sim 10^{-5}$,
$\epsilon_{B}\sim 10^{-5}$ is smaller by $\sim 2-4$ orders of magnitude
compared to the majority of previously reported $\epsilon_{B}$ values (Figure
1), which are $\epsilon_{B}\sim 10^{-3}-10^{-1}$. Assuming $B_{0}\sim
10\mu\mbox{G}$, $\epsilon_{B}\sim\mbox{ few}\times 10^{-5}$ corresponds to
$AF\sim 50$ (Equation 15). Our result of a median $\epsilon_{B}\sim\mbox{
few}\times 10^{-5}$ shows that the majority of the bursts in our X-ray and
optical samples only require a weak amplification beyond shock compression, by
a factor $\sim 50$ or lower.
The near equipartition $\epsilon_{B}\sim 0.01-0.1$ determined near the shock
front by theoretical studies and Particle-In-Cell (PIC) simulations (Medvedev
& Loeb, 1999; Chang et al., 2008; Martins et al., 2009; Keshet et al., 2009;
Lemoine, 2013; Sironi et al., 2013) stands in contrast with our median results
of $\epsilon_{B}\sim\mbox{ few}\times 10^{-5}$. PIC simulations of
relativistic collisionless shocks performed by Chang et al. 2008 and Keshet et
al. 2009 found that the magnetic field generated near the shock front decays
with distance downstream of the shock front. Lemoine 2013 and Lemoine et al.
2013 studied the effects that this decaying magnetic field has on the shock
accelerated electrons radiating afterglow emission downstream of the shock
front. The main effect is that electrons with different Lorentz factors cool
in regions with different magnetic fields, with the higher (lower) energy
photons being emitted by electrons that are closer (further) from the shock
front (Lemoine et al., 2013). Considering the decay of the downstream magnetic
field, Lemoine et al. (2013) modelled the afterglow data of 4 GRBs that have
extended emission at energies $>100\mbox{ MeV}$ (detected by Fermi-LAT) and
also X-ray, optical, and radio data. Their afterglow modelling results for the
X-ray, optical, and radio data found $\epsilon_{B}\sim 10^{-6}-10^{-4}$,
consistent with our results for the median $\epsilon_{B}$ upper limits and
measurements attained from X-ray and optical data.
Our next main result relates to the distribution of $\epsilon_{B}$ values. One
property the $\epsilon_{B}$ values from the literature shared with our optical
$\epsilon_{B}$ measurements is that they both show a wide distribution. The
literature compilation (Figure 1) showed $\epsilon_{B}\sim 10^{-5}-10^{-1}$
and our optical $\epsilon_{B}$ measurement histograms showed an even wider
distribution, ranging from $\epsilon_{B}\sim 10^{-8}-10^{-3}$
($\epsilon_{B}\sim 10^{-9}-10^{-3}$) for a constant density (wind) medium. One
possibility we investigated to explain the wide distribution of $\epsilon_{B}$
values is whether bursts with smaller $\epsilon_{B}$ values are more energetic
than bursts with larger $\epsilon_{B}$ values (Figure 5). Although the bursts
in our optical sample did show the trend that bursts with larger $E$ have a
smaller $\epsilon_{B}$, the correlation was weak, with a correlation
coefficient of $0.62$ and a P-value of $1.2\times 10^{-4}$ (3.8$\sigma$
significance).
Another possibility to explain the wide distribution of $\epsilon_{B}$ relates
to the uncertainty in the environmental parameters ($B_{0}$ and $n$) in the
medium surrounding GRBs. One possibility to explain bursts with values of
$\epsilon_{B}\sim 10^{-5}-10^{-3}$, under the interpretation of weak
amplification beyond shock compression, is that they occurred in environments
with particularly high seed magnetic fields. Since $\epsilon_{B}\propto
B_{0}^{2}$, an increase in $B_{0}$ by an order of magnitude will lead to an
increase in $\epsilon_{B}$ by two orders of magnitude. A weak amplification
beyond shock compression of $AF\sim 50$ (as inferred for the bursts with
$\epsilon_{B}\sim\mbox{ few}\times 10^{-5}$) and a strong $B_{0}\sim
10\mu\mbox{G}-100\mu\mbox{G}$ yields $\epsilon_{B}\sim 10^{-5}-10^{-3}$ (see
Equation 15). We note that for this estimate we assumed $n=1\mbox{ cm}^{-3}$,
but as we discussed in Section 7, the amplification factor has a weak
dependence on the density. This possibility, that the larger values of
$\epsilon_{B}$ may be explained by bursts going off in environments with
particularly strong $B_{0}$, will be discussed further in Barniol Duran 2013
(in preparation).
In addition to many bursts in our optical sample having particularly large
$\epsilon_{B}\sim 10^{-5}-10^{-3}$, there were also some bursts in our optical
sample that have particularly low $\epsilon_{B}$ values. For $p$ determined
from $\alpha_{O}$ and $n=1\mbox{ cm}^{-3}$ (Bottom-right panel of Figure 3),
GRBs 071025 and 071112C have $\epsilon_{B}\sim 10^{-7}$ and GRB 080607 has
$\epsilon_{B}\sim 10^{-8}$. Assuming $n=1\mbox{ cm}^{-3}$ and
$B_{0}=10\mu\mbox{G}$, these bursts with $\epsilon_{B}\sim 10^{-8}-10^{-7}$
are consistent with shock compression of a seed magnetic field of a few 10
$\mu\mbox{G}$ being the only amplification necessary to explain the
observations.
Lastly, we mention that a similar conclusion for radio observations of
supernova remnants (SNRs) was reached in Thompson et al. (2009). Like GRB
afterglow emission, SNR emission arises from a blastwave interacting with the
surrounding medium, but at non-relativistic speeds. Thompson et al. (2009)
found that starburst galaxies have strong ambient magnetic fields $\sim 1$ mG
and they concluded that shock compression of this strong seed magnetic field
is enough to explain the radio emission from SNRs. On the other hand, for
normal spiral galaxies with ambient magnetic fields $\sim 5-10\mu$G, they
concluded that additional amplification beyond shock compression by a factor
$\sim 3-7$ was necessary to explain the radio emission from SNRs.
This work made use of data supplied by the UK Swift Science Data Centre at the
University of Leicester. RS dedicates this work to his friend, Jossue Colato,
thanks Patrick Crumley and Roberto Hernández for helpful discussions, and
Kevin Gately for his help with obtaining the values of $\epsilon_{e}$ and
$\epsilon_{B}$ from the literature. This work has been funded in part by NSF
grant ast-0909110. RBD was supported by an ERC advanced grant (GRB) and by the
I-CORE Program of the PBC and the ISF (grant 1829$/$12).
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## Appendix A $\epsilon_{e}$ And $\epsilon_{B}$ Values From The Literature
In this Appendix, we show a table of the $\epsilon_{e}$ and $\epsilon_{B}$
values we used to make the histograms in Figure 1. For the first GRBs with
high quality afterglow data, different works (e.g. Wijers & Galama 1999;
Chevalier & Li 1999, 2000; Li & Chevalier 2001; Harrison et al. 2001; Yost et
al. 2003) have presented afterglow modeling studies on the same GRBs. In Table
LABEL:epsilon_e_and_epsilon_B_table, for these early afterglow studies, we
show the results from Panaitescu & Kumar 2001a, b; Panaitescu & Kumar 2002
since they have the largest afterglow modeling compilations. It is important
to note that in some cases, different groups determined significantly
different values for the afterglow parameters (e.g. GRB 970508 Wijers & Galama
1999; Chevalier & Li 2000; Frail et al. 2000; Panaitescu & Kumar 2002, GRB
000418 Berger et al. 2001; Panaitescu & Kumar 2002).
Table 6$\epsilon_{e}$ and $\epsilon_{B}$ Values From The Literature GRB | $\epsilon_{e}$ | $\epsilon_{B}$ | Ref.
---|---|---|---
970508 | $0.62$ | $0.10$ | [1]
980329 | $0.12$ | $0.17$ | [2]
980519 | $0.25$ | $(3.5^{+32}_{-2.3})\times 10^{-5}$ | [1]
980703 | $0.14$ | $4.6\times 10^{-4}$ | [3]
990123 | $0.59$ | $(7.4^{+23}_{-5.9})\times 10^{-4}$ | [1]
990510 | $>0.3$ | $6\times 10^{-3}$ | [4]
991208 | $0.32$ | $2.1\times 10^{-2}$ | [1]
991216 | $0.4$ | $2\times 10^{-2}$ | [4]
000301C | $0.4$ | $7\times 10^{-2}$ | [4]
000926 | $0.35$ | $(6.5^{+1.5}_{-1.1})\times 10^{-2}$ | [1]
010222 | $>0.3$ | $2\times 10^{-4}$ | [4]
011211 | $0.22$ | $5.0\times 10^{-4}$ | [5]
020405 | $0.1$ | $0.3$ | [6]
020813 | | $4.0\times 10^{-4}$ | [5]
021004 | $0.21$ | $2\times 10^{-4}$ | [7]
030226 | $0.11$ | $2.5\times 10^{-4}$ | [5]
030329 | $0.16$ | $0.10$ | [8]
050416A | $0.2-0.333$ | $0.072-0.333$ | [9]
050820A | $0.14^{+0.02}_{-0.01}$ | $0.013^{+0.003}_{-0.001}$ | [10]
050904 | $0.02$ | $0.015$ | [11]
051022 | $0.0681^{+0.3951}_{-0.0348}$ | $(8.02^{+28.18}_{-7.17})\times 10^{-3}$ | [12]
051221A | $0.24-0.333$ | $0.12-0.333$ | [13]
060418 | $0.06^{+0.01}_{-0.02}$ | $0.15^{+0.14}_{-0.01}$ | [10]
070125 | $0.27^{+0.03}_{-0.01}$ | $0.0277^{+0.0044}_{-0.0075}$ | [14]
080129 | $0.4$ | $5\times 10^{-2}$ | [15]
080319B | $0.11\pm 0.01$ | $0.33$ | [10]
080928 | $0.165$ | $(2.5^{+16}_{-2.4})\times 10^{-4}$ | [16]
090323 | $0.070^{+0.005}_{-0.005}$ | $0.0089^{+0.0007}_{-0.0018}$ | [17]
090328 | $0.11^{+0.06}_{-0.01}$ | $0.0019^{+0.0004}_{-0.0008}$ | [17]
090423 | $0.28$ | $1.6\times 10^{-4}$ | [18]
Table 6$\epsilon_{e}$ and $\epsilon_{B}$ Values From The Lit. (Continued)
|
arxiv-papers
| 2013-09-12T20:00:03 |
2024-09-04T02:49:50.843720
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Rodolfo Santana, Rodolfo Barniol Duran, Pawan Kumar",
"submitter": "Rodolfo Santana",
"url": "https://arxiv.org/abs/1309.3277"
}
|
1309.3294
|
# A Short Proof that Minimal Sets of Planar Ordinary Differential Equations
are Trivial
Ido Bright Department of Applied Mathematics, University of Washington.
###### Abstract.
We present a short proof, relaying on the divergence theorem, verifying that
minimal sets in the plane are trivial.
‘
## 1\. Introduction
We consider the ordinary differential equation in the plane defined by
(1.1) $\frac{d\mathbf{x}}{dt}=\mathbf{f}(\mathbf{x}),$
where $\mathbf{f}:\mathbb{R}^{2}\rightarrow\mathbb{R}^{2}$ is a locally
Lipschitz function. (Although, we only use uniqueness with respect to initial
conditions of (1.1) and the continuity of $\mathbf{f}$.).
A minimal set is a nonempty closed invariant set, which is minimal with
respect inclusions. A trivial minimal set is a set that is the image if either
a stationary solution or a periodic solution.
We present a new short proof of the following well known result.
###### Theorem 1.
Any minimal set of (1.1) either corresponds to a stationary solution or to the
image of a periodic solution, namely, all minimal sets are trivial.
The text-book proof of this theorem relays on the Poincaré–Bendixson theorem,
and employs dynamical arguments. The proof presented in this paper relays on a
different argument, relaying on a property of the velocity of Jordan curves.
This idea was introduced in [1] and further developed in [2, 3].
## 2\. Proof of Main Result
In the proof of the main result we use the following notation. The
2-dimensional euclidean space is denoted by $\mathbb{R}^{2}$, and the norm of
a vector $\mathbf{y}\in\mathbb{R}^{2}$ is denoted by
$\left|\mathbf{y}\right|$. The open ball in $\mathbb{R}^{2}$, centered at
$\mathbf{y}$ with radius $r$, is denoted by $B\left(\mathbf{y},r\right)$. The
closure of an open set $O\subset\mathbb{R}^{2}$ is denoted by $\bar{O}$, its
boundary by $\partial O$, and its exterior normal and tangent vector at the
point $\mathbf{y}\in\partial O$ are denoted by $\mathbf{N}_{\partial
O}\left(\mathbf{y}\right)$ and $\mathbf{T}_{\partial
O}\left(\mathbf{y}\right)$, respectively.
We shall use the following results that are well known in the smooth case.
###### Lemma 2.
Let $O\subset\mathbb{R}^{2}$ be a bounded open set with rectifiable boundary.
Then
$\mathbf{v}=\int_{\partial O}\mathbf{N}_{\partial
O}\left(\mathbf{y}\right)d\mathbf{y}=\mathbf{0}\in\mathbb{R}^{2}.$
###### Proof.
Assume in contradiction that $\mathbf{v}\neq\mathbf{0}$, and set
$\mathbf{g}:\mathbb{R}^{2}\rightarrow\mathbb{R}^{2}$ by
$\mathbf{g}\equiv\frac{\mathbf{v}}{\left|\mathbf{v}\right|}$. The divergence
theorem for sets of finite perimeter (see, e.g., [4]) implies that
$\left|\mathbf{v}\right|=\frac{\mathbf{v}}{\left|\mathbf{v}\right|}\mathbf{v}=\left|\int_{\partial
O}\mathbf{g}\left(\mathbf{y}\right)\mathbf{N}_{\partial
O}\left(\mathbf{y}\right)d\mathbf{y}\right|=\left|\int_{O}\nabla\cdot\mathbf{g}\left(\mathbf{y}\right)d\mathbf{y}\right|=0,$
in contradiction. ∎
The following lemma appears in [5].
###### Lemma 3.
Suppose $I\subset\mathbb{R}$ is a bounded interval and
$g:I\rightarrow\mathbb{R}$ is a Lipschitz function. Then for almost every
$r\in\mathbb{R}$ the set $g^{-1}\left(r\right)=\left\\{t\in
I|g\left(t\right)=r\right\\}$ is finite.
To prove the main theorem, let us now fix a minimal set
$\Omega\subset\mathbb{R}^{2}$ and a solution
$\mathbf{x}^{*}\left(\cdot\right)$ of (1.1), defined on $[0,\infty)$, with
trajectory contained in $\Omega$.
We shall also use the following well known fact.
###### Lemma 4.
For every $\mathbf{y_{0}}\in\Omega$ and $\delta>0$ there exists $t>s$ such
that $\left|\mathbf{x}^{*}\left(t\right)-\mathbf{y_{0}}\right|<\delta$.
###### Proof.
Otherwise, suppose that the lamma does not hold for some
$\mathbf{y_{0}},\delta$ and $s$. Then the curve
$\mathbf{y}^{*}\left(t\right)=\mathbf{x}^{*}\left(s+t\right)$ is a solution of
(1.1) with trajectory contained in $\Omega\backslash
B\left(\mathbf{y_{0}},\delta\right)$ for a suitable $\delta>0$, in
contradiction to the minimality of $\Omega$. ∎
If $\Omega$ is not a singleton we choose $D>0$ such that $\Omega\backslash
B\left(\mathbf{x}^{*}\left(0\right),3D\right)\neq\emptyset$ and apply the
following construction:
###### Construction 5.
Set $\delta_{0}=D$ and $t_{0}$ as the first time point where
$\mathbf{x}^{*}\left(\cdot\right)$ meets $\partial
B\left(\mathbf{x}^{*}\left(0\right),\delta_{0}\right)$. For $i=1,2,\dots$ do
the following:
1. (1)
Choose $\delta_{i}<\delta_{i-1}/2$ small enough, such that
$\left|\mathbf{x}^{*}\left(0\right)-\mathbf{x}^{*}\left(t\right)\right|>\delta_{i}$
for all $t\in\left[t_{0},t_{i-1}\right]$.
2. (2)
Set $t_{i}$ as the first time point after $t_{0}$ where the curve
$\mathbf{x}^{*}\left(\cdot\right)$ meets $\partial
B\left(\mathbf{x}^{*}\left(0\right),\delta_{i}\right)$. (Here we use Lemma 4).
3. (3)
Starting from $\mathbf{x}^{*}\left(t_{i}\right)$ follow the line connecting it
to $\mathbf{x}^{*}\left(0\right)$, until first meeting a point in
$\mathbf{x}^{*}\left(\left[0,t_{0}\right]\right)$. Let
$\mathbf{x}^{*}\left(s_{i}\right)$ be this point.
4. (4)
Let $\gamma_{i}$ be the parametrized Jordan curve obtained by following the
curve $\mathbf{x}^{*}\left(\cdot\right)$ in the interval
$\left[s_{i},t_{i}\right]$ and then the line connecting its endpoints, with
velocity of norm $1$.
###### Lemma 6.
If $t_{i}\rightarrow t^{*}$ then
$\mathbf{x}^{*}\left(0\right)=\mathbf{x}^{*}\left(t^{*}\right)$, and
$\mathbf{x}^{*}\left(\cdot\right)$ is a periodic solution with image $\Omega$.
###### Proof.
According to our construction
$\left|\mathbf{x}^{*}\left(0\right)-\mathbf{x}^{*}\left(t_{i}\right)\right|=\delta_{i}<2^{-i}D$
for every $i$. Hence, by continuity
$\mathbf{x}^{*}\left(t^{*}\right)=\mathbf{x}^{*}\left(0\right)$, and
$\mathbf{x}^{*}\left(\cdot\right)$ is periodic. By the minimality of $\Omega$,
the image of $\mathbf{x}^{*}\left(\cdot\right)$ is $\Omega$. ∎
###### Proof of Theorem 1.
Clearly, $\Omega$ is a singleton if and only if it contains a point
$\mathbf{y}\in\Omega$ such that
$\mathbf{f}\left(\mathbf{y}\right)=\mathbf{0}$. In this case and when the
condition of Lemma 6 holds, we are done. Thus, we assume that $\mathbf{f}$
does not vanish in $\Omega$ and that $t_{i}\rightarrow\infty$.
Fix $\mathbf{y}_{0}\in\Omega$ such that
$\left|\mathbf{y}_{0}-\mathbf{x}^{*}\left(0\right)\right|>2D$. Using Lemma 3
we fix an arbitrary small ball $B=B\left(\mathbf{y}_{0},r_{0}\right)$, such
that $r_{0}<D$, and that $\left\\{0\leq t\leq
s|\left|\mathbf{x}^{*}\left(t\right)-\mathbf{y_{0}}\right|=r_{0}\right\\}$ is
finite for every $s>0$. Note that this implies that for every $i$ the Jordan
curve $\gamma_{i}$ intersects $\partial B$ at a finite number of points, and
that the portion of $\gamma_{i}$ in $B$ corresponds to the trajectory
$\mathbf{x}^{*}\left(\cdot\right)$.
For every $i$ we denote the interior of $\gamma_{i}$ by $O_{i}$, and, using
the identity
$\partial\left(O_{i}\cap B\right)\subset\left(\partial O_{i}\cap
B\right)\cup\left(O_{i}\cap\partial B\right)\cup\left(\partial
O_{i}\cap\partial B\right),$
we obtain, by Lemma 2, that
$\mathbf{0}=\int_{\partial\left(O_{i}\cap
B\right)}\mathbf{N}_{\partial\left(O_{i}\cap
B\right)}\left(\mathbf{y}\right)d\mathbf{y}=\int_{\partial O_{i}\cap
B}\mathbf{N}_{\partial
O_{i}}\left(\mathbf{y}\right)d\mathbf{y}-\int_{O_{i}\cap\partial
B}\mathbf{N}_{\partial B}\left(\mathbf{y}\right)d\mathbf{y},$
since $\partial O_{i}\cap\partial B$ has zero measure. This bounds
(2.1) $\left|\int_{\partial O_{i}\cap B}\mathbf{N}_{\partial
O_{i}}\left(\mathbf{y}\right)d\mathbf{y}\right|=\left|\int_{O_{i}\cap\partial
B}\mathbf{N}_{\partial
B}\left(\mathbf{y}\right)d\mathbf{y}\right|\leq\left|\int_{O_{i}\cap\partial
B}\left|\mathbf{N}_{\partial
B}\left(\mathbf{y}\right)\right|d\mathbf{y}\right|\leq 2\pi r_{0}.$
For each $i$ the set $\partial O_{i}\cap B$ contains a finite number of arcs,
and applying a change of variable it is easy to see that
$\int_{\partial O_{i}\cap B}\mathbf{N}_{\partial
O_{i}}\left(\mathbf{y}\right)d\mathbf{y}=P_{i}\int_{\partial O_{i}\cap
B}\mathbf{T}_{\partial
O_{i}}\left(\mathbf{y}\right)d\mathbf{y}=P_{i}\int_{\left\\{t\leq
t_{i}|\mathbf{x}^{*}\left(t\right)\in
B\right\\}}\frac{d}{dt}\mathbf{x}^{*}\left(t\right)dt,$
where $\mathbf{T}_{\partial O_{i}}$ is chosen to agree with the direction of
$\gamma_{i}$, and $P_{i}$ is a $\frac{\pi}{2}$-rotation matrix. Here we use
the fact that the portion of $\gamma_{i}$ in $B$ corresponds to the original
trajectory $\mathbf{x}^{*}\left(\cdot\right)$.
Combined with (2.1) we conclude that for every $i$
$\left|\int_{\left\\{t\leq t_{i}|\mathbf{x}^{*}\left(t\right)\in
B\right\\}}\mathbf{f}\left(\mathbf{x}^{*}\left(t\right)\right)dt\right|\leq
2\pi r_{0}.$
The minimality of $\Omega$ and Lemma 4 implies that the set
$\left\\{t|\mathbf{x}^{*}\left(t\right)\in B\right\\}$ has infinite measure.
This implies that $\mathbf{0}$ is contained in the convex hull of
$\left\\{\mathbf{f}\left(\mathbf{y}\right)|\mathbf{y}\in\bar{B}\right\\}$. The
radius $r_{0}$ can be chosen arbitrary small, thus, the continuity of
$\mathbf{f}$ implies that $\mathbf{f}\left(\mathbf{y}_{0}\right)=\mathbf{0}$,
in contradiction.
## References
* [1] Z. Artsein and I. Bright. Periodic optimization suffices for infinite horizon planar optimal control. SIAM Journal on Control and Optimization, 48:4963–4986, 2010.
* [2] I. Bright. A reduction of topological infinite-horizon optimization to periodic optimization in a class of compact 2-manifolds. Journal of Mathematical Analysis and Applications, 394(1):84 – 101, 2012.
* [3] I. Bright and J.M. Lee. Bounds for flux integrals over regular domains. Submitted.
* [4] W.F. Pfeffer. The Divergence Theorem and the Sets of Finite Perimeter, volume 303\. CRC Press, 2012.
* [5] J.C. Wells. A note on a corollary of sardś theorem. Proceedings of the American Mathematical Society, 48(2):513–514, 1975.
∎
|
arxiv-papers
| 2013-09-12T20:36:26 |
2024-09-04T02:49:50.862951
|
{
"license": "Public Domain",
"authors": "Ido Bright",
"submitter": "Ido Bright",
"url": "https://arxiv.org/abs/1309.3294"
}
|
1309.3524
|
Gakuto International Series
Mathematical Sciences and Applications Vol.XX (2011)
International Symposium on Computational Science 2011
pp.XXXXX
Stochastic moving particle semi-implicit for inviscid fluid wave simulation
Christian Fredy Naa
Institute of Science and Engineering, School of Mathematics and Physics,
Kanazawa University
Faculty of Mathematics and Natural Science, Institut Teknologi Bandung
Seiro Omata
Institute of Science and Engineering, School of Mathematics and Physics,
Kanazawa University
Masaki Kazama
Fujitsu Limited
Abstract. The present paper introduces stochastic velocity as improvement for
moving particle semi-implicit (MPS) method. This improvement is to overcome
energy loss caused by numerical dissipation in the basic MPS that brings about
rapid decay of waves. Stochastic velocity is added in the explicit step of the
basic MPS method. MPS with stochastic improvement is compared with the basic
method in the case of linear water waves, in particular dam break problem and
standing wave in a rectangular tank. Surface detection and curve fitting are
used to analyze the parameters of wave on the standing wave case. The surface
detection and curved fitting was efficient to determine parameters of the wave
and it was found that the stochastic improvement made the waves survived
longer than in the basis method.
——————————————————————– ———-
Received xxxxxxxxxx, 2011.
This work is supported by xxxxxxxxxxxxxxxxxxxx.
AMS Subject Classification xxxxxxxxxx, xxxxxxxxxx
## 1 Introduction
Moving particle semi-implicit (MPS) [1] method is a particle method for
simulating incompressible fluids. MPS method was used for example to analyze
breaking waves [2], droplet breakup behavior [3] and to predict wave impact
pressure [4]. Since MPS method is based on Lagrangian system, computational
grids are not necessary. Governing equations are discretized based on particle
interaction models representing density, gradient, Laplace operator and free
surface.
However, MPS has a weakness in the energy conservation. The waves calculated
by the basic MPS decay rapidly since the mechanical energy is not fully
conserved. This lack of conservation of energy is caused by numerical
dissipation. A work has been introduced to recover this weakness: Suzuki [5]
introduced Hamiltonian moving particle semi-implicit (HMPS). Even then, the
mechanical energy was not fully conserved, but the HMPS was able to make the
waves survive longer.
The objective of this paper is to overcome the loss of energy in the basic MPS
using stochastic modification of velocity. The purpose of this stochastic
concept is to add extra kinetic energy to the particles so that the kinetic
energy gained by the stochastic velocity recovers the loss of energy caused by
numerical dissipation.
The continuity equation and Euler’s equation are used as governing equation.
Stochastic and basic MPS were compared in the cases of dam break problem and
standing wave in a rectangular tank.
Surface detection algorithm is used to determine amplitude in the standing
wave case. This kind of algorithm is used because it is hard to judge the
important parameters of the wave only from the distribution of particles.
The paper is organized in the following way. A brief explanation of standard
MPS method is presented in second section. The stochastic improvement of
velocity is introduced in third section. The fourth section describes the
surface detection algorithm. Finally, the improved MPS method compared with
the standard MPS method in the case of dam break problem and standing wave in
a rectangular tank.
## 2 Standard MPS method
In this section, the MPS method is briefly explained based on description
which provided by Koshizuka [1]. In the MPS method, the fluid is modeled using
interaction of particles according to equations of motion. Governing equations
for inviscid fluid motion are continuity equation and Euler’s equation:
$\frac{1}{\rho}\frac{D\rho}{Dt}+\nabla\cdot{\bf{u}}=0$ (1)
$\frac{D\bf{{u}}}{Dt}=-\frac{1}{\rho}\nabla P+{\bf{g}},$ (2)
where ${\bf{u}}$ denotes particle velocity vector, $t$ denotes time, $\rho$
denotes density, $P$ denotes pressure and ${\bf{g}}$ denotes gravity
acceleration vector.
It should be noted that Eq. (1) is written in the form of a compressible flow.
In the MPS method, incompressibility is enforced by the way setting
$\frac{D\rho}{Dt}=0$ at each particle at each calculation time step. According
to [4] the left hand side of Eq. (2) denotes the material derivative
$\frac{D\bf{{u}}}{Dt}$ involving the advection term. In the particle methods,
including the MPS method, the advection term is automatically calculated
through the tracking of particle motion; hence, the numerical diffusion
arising from the successive interpolation of the advection function in
Eulerian grid based methods is controlled without the need for a sophisticated
algorithm.
The main idea of MPS method is to divide Eq. (2) into two parts to calculate
the change of velocity, as follows
$\left(\frac{d{\bf{u}}}{{dt}}\right)^{\text{explicit}}={\bf{g}},$ (3)
$\left(\frac{d{\bf{u}}}{dt}\right)^{\text{implicit}}=-\frac{1}{\rho}\nabla P.$
(4)
Particle interaction is described in terms of weight function. Weight function
$w(r)$ in MPS method is defined as [1]
$w(r)=\begin{cases}\frac{r_{e}}{r}-1&\text{if }r<r_{e}\\\ 0&\text{if }r\geq
r_{e}.\end{cases}$ (5)
Here, $r$ will have the meaning of distance between particles and $r_{e}$ is
the _cut off_ distance. In this paper, $r_{e}$ equals to $2.1l_{0}$, where
$l_{0}$ is the initial distance between particles.
The particle number density for particle $i$ ($n_{i}$) is calculated by
$n_{i}=\sum_{j\neq i}^{N}w(r),$ (6)
where $N$ denotes the total number of particles.
The Laplace operator of a scalar quantity $\phi$ for particle $i$ is evaluated
using
$\left<\nabla^{2}\phi\right>_{i}=\frac{2d}{\lambda n_{0}}\sum_{j\neq
i}^{N}(\phi_{j}-\phi_{i})w(|{\bf{r}}_{j}-{\bf{r}}_{i}|).$ (7)
Here, ${\bf{r_{i}}}$ is position vector of particle $i$ and $d$ is the space
dimension. In this paper, the constant $n_{0}$ is defined as the maximum value
of particle number density
$n_{0}=\max{n_{i}}.$ (8)
The parameter $\lambda$ is defined as [4]
$\lambda=\frac{\sum_{i\neq
j}^{N}|{\bf{r}}_{j}-{\bf{r}}_{i}|^{2}w(|{\bf{r}}_{j}-{\bf{r}}_{i}|)}{\sum_{i\neq
j}^{N}w(|{\bf{r}}_{j}-{\bf{r}}_{i}|)}.$ (9)
The gradient of pressure $\nabla P$ for particle $i$ is defined as
$\left<\nabla P\right>_{i}=\frac{d}{n_{0}}\sum_{j\neq
i}^{N}\frac{P_{j}-\hat{P}_{i}}{|{\bf{r}}_{j}-{\bf{r}}_{i}|^{2}}w(|{\bf{r}}_{j}-{\bf{r}}_{i}|)({\bf{r}}_{j}-{\bf{r}}_{i}),$
(10)
where $\hat{P}_{i}$ denotes the minimum pressure among particles within
certain cut-off distance.
For modeling the incompressibility, the number of densities $n^{*}$ that are
calculated at the end of explicit step deviate from the constant number of
density $n_{0}$; hence, a second corrective process is required to adjust the
number of densities to initial values prior to the time step. In the implicit
step, the intermediate particle velocities are updated implicitply through
solving the Poisson Pressure Equation (PPE) derived as [1]
$\left<\nabla^{2}P_{k+1}\right>_{i}=\frac{\rho}{\Delta
t^{2}}\frac{n_{0}-(n_{k}^{*})_{i}}{n_{0}},$ (11)
where $\Delta t$ denotes calculation time step and $k$ denotes the step of
calculation.
To set the Dirichlet boundary condition for the Poisson’s equation of
pressure, particles satisfying [3]
$n_{i}<\beta n_{0},$ (12)
are judged as surface particles and their pressure is fixed to zero or to
atmospheric pressure value. Here, $\beta$ is a value between $0.8$ to $0.98$.
The time step is important for numerical stability. According to [6] the time
step should follow the CFL condition
$\Delta t\leq 0.2\frac{l_{0}}{u_{\text{max}}},$ (13)
here, $u_{\text{max}}$ is the maximum velocity among the particles.
## 3 Modification of MPS method
As described before that the standard MPS method has a weakness in the energy
conservation. The mechanical energy is not fully conserved caused by numerical
dissipation. To overcome this problem, the stochastic improvement of velocity
is introduced. The stochastic velocity is added in standard MPS method.
To add such stochastic velocity, after the particle’s velocity in the explicit
step (Eq. (3)) is calculated, the stochastic velocity ${\bf{u}}^{\text{stoc}}$
is added. The direction of stochastic velocity ${\bf{u}}^{\text{stoc}}$ should
be determined so that it does not reduce the particle velocity in the
particular time step. To generate such direction, first the unit vector of
particle’s velocity u is calculated by
$\begin{cases}\hat{u}_{x}=\frac{u_{x}}{|{\bf{u}}|}\\\
\hat{u}_{y}=\frac{u_{y}}{|{\bf{u}}|},\end{cases}$ (14)
where $|{\bf{u}}|$ is the length of the vector. Next, random angle $\theta$ is
generated from the interval ($-\pi/2$,$\pi/2$) by choosing a random number $r$
from a uniform distribution $U(0,1)$:
$\theta=r\left(\pi/2-(-\pi/2)\right)-\pi/2.$ (15)
Then, rotating the vector $\hat{u}_{x}$ and $\hat{u}_{y}$ by the angle
$\theta$ counter-clockwise about the origin, the unit vector of the stochastic
velocity is determined by
$\left(\begin{array}[]{cc}u_{x}^{\text{stoc}}\\\
u_{y}^{\text{stoc}}\end{array}\right)=\left(\begin{array}[]{cc}\cos\theta&-\sin\theta\\\
\sin\theta&\cos\theta\end{array}\right)\left(\begin{array}[]{cc}{\hat{u}_{x}}\\\
{\hat{u}_{y}}\end{array}\right)$ (16)
Finally, the unit vector determined by Eq. (16) is multiplied by the magnitude
of the stochastic velocity $m$ to get the final vector form
${\bf{u}}^{\text{res}}={\bf{u}}+m{\bf{u}}^{\text{stoc}}.$ (17)
Here, u is the velocity calculated by basic MPS method (after the explicit
step). ${\bf{u}}^{\text{res}}$ is the resultant velocity due to addition
${\bf{u}}$ and ${\bf{u}}^{stoc}$ This process is illustrated in Fig. 1.
Figure 1: Stochastic velocity
The magnitude of the stochastic velocity $m$ is also given by choosing a
random number $r$ from a uniform distribution $U(0,1)$:
$m=r\left(\text{max}-\text{min}\right)+\text{min},$ (18)
where min is the minimum range of the magnitude which is defined to be $0$,
max is the maximum range which defined as
$\text{max}=\alpha u_{\text{loss}}.$ (19)
Here, $\alpha$ is a positive constant which determines the strength of the
stochastic velocity, while the term $u_{\text{loss}}$ is defined as the amount
of velocity that is lost due to energy difference on subsequent time steps.
The value of $u_{\text{loss}}$ is determined based on the idea that the energy
loss on two neighboring time steps corresponds to the loss of kinetic energy
$\Delta EK$
$\Delta EK=E_{\text{total}}^{k-1}-E_{\text{total}}^{k},$ (20)
where $E_{\text{total}}$ denotes total energy and $k$ denotes step number.
Since kinetic energy per unit mass is given by $EK=\frac{{\bf{|u|}}^{2}}{2}$,
then
$u_{\text{loss}}=\sqrt{2\Delta EK}.$ (21)
The addition of stochastic velocity follows several constraints:
1. 1.
From experience of several simulations, the constant $\alpha$ is usually much
less than $0.1$. If its more than $0.1$, the particles will get extra high
velocity that leads to excessive energy making the particles jump off.
2. 2.
It may happen that $\Delta EK$ in Eq. (20) equals to zero or even to a
negative number (which means the energy is larger than in the previous step).
In this case, the explicit step will be calculated without stochastic
velocity.
The complete algorithm of stochastic MPS in some step $k$ is described below:
Algorithm 1 Stochastic MPS method algorithm
1:Total energy $E^{k-1}_{total}$
2:Calculate total energy $E^{k}_{\text{total}}$.
3:Determine direction and magnitude ${\bf{u}}^{\text{stoc}}$
4:if $E^{j-1}_{\text{total}}\leq E^{k}_{\text{total}}$ then
5: ${\bf{u}}^{\text{res}}={\bf{u}}+m{\bf{u}}^{\text{stoc}}$
6:else
7: ${\bf{u}}^{\text{res}}={\bf{u}}$
8:end if
9:$E^{k-1}_{\text{total}}=E^{k}_{\text{total}}$.
10:check termination.
This paper concerns the addition of stochastic velocity in 2D case; similar
idea can be adopted for 3D case. Since this stochastic concept is a new
improvement in MPS method, it is not yet verified by any other research.
## 4 Surface Detection
In order to compare the result from the simulations of basic MPS method and
the stochastic improvement, a surface detection algorithm is needed. It is
hard to judge the important parameters of the wave only from the distribution
of particles.
The algorithm begins with finding the upper left and right particles. Then the
algorithm starts to trace the surface particles from the left to the right.
From the left particle, the vertical reference axis is used to detect the
second surface particle.
Figure 2: Detection by using vertical imaginary axis. The solid circles denote
the surface solution.
The mechanism to detect the second surface particles is shown in Fig. 2. The
radius of search area is bounded by constant $r_{e}$ which is the same as the
cut-off radius in the above described theory of MPS. Angle from the vertical
reference axis to each particle within the search radius is measured. Once the
angles are measured, they are converted into quadrant value. The second
surface particle is then determined as the particle that has the minimum angle
with the vertical reference axis. After the second surface particle is
determined, this surface particle is called the _current point_ and it is used
to get the third surface particle.
Figure 3: Detection using dot product relation. The solid circles denote the
surface solution.
The third surface particle and so on are determined using different algorithm.
This algorithm is shown in Fig. 3, the surface particle detected before the
current point is used (in this case: the third surface particle is determined
using the first and the second particles).
First, the reference vector that connects the current point and the previous
surface particle is measured. Then the normal vector perpendicular to the
reference vector which direct out from the particle distribution is determined
using geometry transformation matrix. Candidate vectors from current point to
all the particles inside the search area are also measured.
Using two dot products relation (the product of between the reference vector
and the candidate vectors, and the product of normal vector and the candidate
vectors), the angles $\theta$ from the reference vector to the candidate
vectors as shown in Fig. 3 can be determined. The next surface particle is the
one corresponding to the smallest angle $\theta$.
Figure 4: Surface detection algorithm
Once the next surface particle is determined, this particle becomes the
current point and the dot product algorithm is used to get the next surface
particle. This mechanism is repeated until the upper right particle is
detected. The flow chat of the algorithm is shown in Fig. 4.
In some simulations, particles jump off. The algorithm should consider a
constraint that would exclude this kind of particles because these particles
are actually not part of the wave. The value of particle number density is
used as a constraint. Namely If the particle number density is less than
$0.4n_{0}$, the algorithm will exclude these particles from the surface
detection process.
## 5 Test Cases
### 5.1 Dam break problem
Two-dimensional dam break problem were simulated. This kind of simulation has
been the most common test case in fluid dynamics.
#### 5.1.1 Conditions of computation
The computation domain was set to be $((x,y)|0\leq x\leq l,0\leq y\leq l)$,
where $l$ was taken to be $2$ m. The initial water column was set to be
$((x,y)|0\leq x\leq 0.5,0\leq y\leq 1)$. Rigid boundary condition was applied.
The initial velocity and pressure was set to zero for each particle. The
distance between particles was set to be $0.02$ m and $2.171$ particles were
used in the computation while the simulation is taken until 15 seconds. The
initial configuration of the particles is illustrated in Fig. 5.
Figure 5: Initial condition of dam break problem
#### 5.1.2 Results
In this implementation, the basic MPS and stochastic MPS were simulated.
Stochastic MPS used several $\alpha$ values (see (19)). Namely,
$\alpha=0.003$, $\alpha=0.001$ and $\alpha=0.0009$ were simulated.
The results of the simulation are shown in Fig. 6. Since the time step on each
simulation is different, then the comparison was taken on the slightly
different time and the main focus of the comparison is the time when wave
disappeared . As shown on Fig. 6, at time $12.76s$ It is shown that the wave
from the basic MPS already disappeared while the stochastic MPS still survived
until several cycle. Thus it can be concluded that the wave which performed by
stochastic improvement stronger and survived longer than the basic MPS.
$\begin{array}[]{cc}\includegraphics[scale={0.32}]{dam0_2eps}&\includegraphics[scale={0.32}]{dam1_2eps}\\\
\includegraphics[scale={0.32}]{dam0_6eps}&\includegraphics[scale={0.32}]{dam1_6eps}\\\
\includegraphics[scale={0.32}]{dam0_10eps}&\includegraphics[scale={0.32}]{dam1_10eps}\\\
\end{array}$
Figure 6: Comparison of dam break results between the stochastic MPS with
$\alpha=0.003$ (left) and the basic MPS (right)
The total energy density with respect to time is plotted in Fig. 7. It is
shows that the energy performed by basic MPS has lower state than the
stochastic. The oscilation of stochastic energy density is caused by the
current algorithm that perform the addition of stochastic velocity only if the
total energy on the particular step has lower value than the step before.
Since the whole energy system will decrease, the oscilation will occur because
of the current stochastic velocity algorithm.
Figure 7: Comparison between basic MPS and the stochastic MPS, $\alpha=0.003$,
$\alpha=0.001$ and $\alpha=0.0009$ according to total energy density
### 5.2 Standing wave in a rectangular tank
Two-dimensional standing waves in a rectangular tank were simulated. This kind
of simulation has been studied by Suzuki [5] using Hamiltonian moving particle
semi-implicit (HMPS). The results are compared with the analytical solution
according to Wu and Taylor [7] based on the water elevation at the center of
the tank. In this paper, different approach using surface detection and least
square curve fitting is used to determine the amplitude of the wave.
#### 5.2.1 Conditions of computation
The computation domain was set to be $((x,y)|0\leq x\leq l,0\leq y\leq l)$,
where $l$ was taken to be equal to the wavelength. Here the wavelength
$\lambda$ was $1$ m. The depth of water $h$ was $\lambda/3$. The initial
configuration of the free surface is given by
$y_{0}(x)=A\cos[k(x+l/2)].$ (22)
Here $y_{0}$ is the initial surface displacement, $A=0.07\lambda$ is the
amplitude and $k$ is the wave number which defined as $2\pi/\lambda$. Periodic
boundary conditions were applied, while at the bottom mirror boundary
condition was used. The initial velocity and pressure was set to zero for each
particle. The distance between particles was set to be $0.02$ m and 834
particles were used in the computation while the simulation is taken until $6$
seconds. The initial configuration of the particles is illustrated in Fig. 8.
Figure 8: Initial condition of standing wave in a rectangular tank
Surface detection was applied to the solution and the obtained surface was
fitted using least square curve fitting. _GNUPLOT_ curve fitting is used. In
this case, since the surface seems closed to quadratic function then the data
fitted using simple quadratic function
$y(x)=ax^{2}+bx+c,$ (23)
where constant $a$, $b$ and $c$ are to determined.
From the result of the least square curve fitting, the amplitude of the wave
is calculated on each time step. From the quadratic function (Eq. 23), the
maximum point ($x_{\text{max}},y_{\text{max}}$) is calculated by setting
$\frac{dy}{dx}=0$ which results in
$\begin{cases}x_{\text{max}}=\frac{-b}{2a}\\\
y_{\text{max}}=ax_{\text{max}}^{2}+bx_{\text{max}}+c.\end{cases}$ (24)
To estimated the quality of the fitting, average relative error
$\bar{\epsilon}$ is used
$\bar{\epsilon}=\frac{1}{N}\sum_{i=0}^{N}\left(\frac{|y_{\text{eq}}-y_{\text{est}}|}{y_{\text{eq}}}\right)\times
100\%.$ (25)
Here, $N$ denotes the number of data, $y_{\text{eq}}$ denotes the fitted value
according to (23), while $y_{\text{est}}$ is the $y$ value of the surface
detection result. Figure 9 shows the result of surface detection and error
calculation.
Figure 9: Example of surface detection and curve fitting result
#### 5.2.2 Results
In this implementation, the basic MPS and stochastic MPS were simulated.
Stochastic MPS used several $\alpha$ values (see (19)). Namely, $\alpha=0.02$,
$\alpha=0.03$ and $\alpha=0.04$ were simulated.
According to the relative error which was less than $4.5\%$, it can be
concluded that the curve fitting worked well.
Figure 10: Comparison between basic MPS, stochastic MPS with $\alpha=0.04$ and
the analytic solution according to wave height (water elevation) in the center
of the tank.
The wave height evolution with respect to time is plotted in Fig. 10. It shows
that the waves obtained from stochastic MPS survived longer than the basic
MPS. The wave from basic MPS decayed rapidly in time, and disappeared already
from $t=3.5$ seconds. In comparison with the analytical solution [7], the wave
height from stochastic MPS was higher and nearer to the analytical solution.
Ideally in qualitative analysis, if the system has perfect mechanical energy
conservation, the amplitude of the wave will not decrease. But since the MPS
method suffers from the numerical dissipation, the wave will rapidly decay.
Using the stochastic velocity, the rapid decay of the wave could be prevented.
Figure 11: Comparison between basic MPS and the stochastic MPS using
$\alpha=0.04$ according to total energy density
The energy density with respect to time is plotted in Fig. 11. It is shown
that energy of the basic MPS decayed faster in time compared with the
stochastic result (with $\alpha=0.04$).
However, the stochastic MPS has a drawback that the particles moved to one
direction after the wave disappeared. This is caused by the random movement of
the particles plus extra velocity added by the stochastic term. The
development of stochastic MPS without this drawback is left for further study.
## 6 Conclusion
This paper introduces an improvement of moving particle semi-implicit method
through the addition of stochastic velocity. Stochastic velocity is added
after the explicit step of basic MPS method. The cases that studied were the
dam break problem and standing wave in a rectangular tank. Surface detection
and least square curve fitting were used to analyze the waves on the standing
wave case.
On the dam break problem, it was showed that the wave that performed by the
stochastic survived longer than the basic MPS. However, the energy density
from the stochastic MPS suffers some oscilation due to the stochastic
algorithm.
On standing wave case, it was showed that the wave with the stochastic MPS
survived longer than the basic MPS. In comparison with the analytical solution
[7], the wave height from stochastic MPS was higher and nearer to the
analytical solution even it still far from the ideal system in which the
mechanical energy is fully conserved.
However, on the standing wave case the stochastic MPS has a drawback that the
particles moved to one direction after the wave disappeared. This is caused by
the random movement of the particles plus extra velocity added by the
stochastic term. The development of stochastic MPS without this drawback is
left for further study.
## References
* [1] Koshizuka S., Oka Y.: _Moving particle semi implicit method for fragmentation of incompressible fluid_. Nuclear Science and Engineering, 123(3):421-434, 1996.
* [2] Koshizuka S., Nobe A. and Oka Y. : _Numerical analysis of breaking waves using the moving particle semi-implicit method_. International Journal of Numerical Methods in Fluids, 26: 751-769, 1998.
* [3] Nomura et. al: _Numerical Analysis of Droplet Breakup Behaviour using Particle Method_. Journal of Nuclear Science and Technology, 38 (12):1057-1064, 2001.
* [4] Khayyer A., Gotoh H. : _Modified Moving Particle Semi-implicit method for the prediction of 2D wave impact pressure_. Coastal Engineering, 56: 419-440, 2009.
* [5] Suzuki Y., Koshizuka S., Oka Y. : _Hamiltonian moving-particle semi implicit (HMPS) method for incompressible fluid flows_. Computer methods in applied mechanics and engineering, 196: 2876-2894, 2007.
* [6] Zhang et.all, _An improved MPS method for numerical simulations of convective heat transfer problem_. International Journal for Numerical Method in Fluids, 51:31-47, 2005.
* [7] Wu G.X., Taylor R. E. : _Finite element analysis of two-dimensional non-linear transient water waves_. Appl. Ocean Res., 16:363-372, 1994.
* [8]
|
arxiv-papers
| 2013-09-13T17:44:23 |
2024-09-04T02:49:50.883048
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Christian Fredy Naa, Seiro Omata, Masaki Kazama",
"submitter": "Christian Naa",
"url": "https://arxiv.org/abs/1309.3524"
}
|
1309.3632
|
3D Simulation of Dam-break effect on a Solid Wall using
Smoothed Particle Hydrodynamic
Suprijadia,b, F. Faizalb, C.F. Naaa and A.Trisnawanb
aDepartment of Physics, bDepartment of Computational Sciences
Faculty of Mathematics and Natural Science, Institut Teknologi Bandung
Jl. Ganesha No.10, Bandung 40132, Indonesia
email : [email protected]
Abstract. Dam is built for water supply, water flow or flooding control and
electricity energy storage, but in other hand, dam is one of the most
dangerous natural disaster in many countries including in Indonesia. The
impact of dam break in neighbour area and is huge and many flooding in remote
area, as happen in Dam Situ Gintung in Tangerang (close to Jakarta) in 2009.
Smoothed Particle Hydrodynamics (SPH), is one of numerical method based on
Lagrangian grid which is applied in astrophysical simulation may be used to
solve the simulation on dam break effect. The development of SPH methods
become alternative methods to solving Navier Stokes equation, which is main
key in fluid dynamic simulation. In this paper, SPH is developed for
supporting solid particles in use for 3D dam break effect (3D-DBE) simulation.
Solid particle have been treated same as fluid particles with additional
calculation for converting gained position became translation and rotation of
solid object in a whole body. With this capability, the result of 3D-DBE
simulation has been varies and interesting. The goals of this simulation is
for analyse fluid and solid particle interaction by using two different
scenario. The first scenario relation between height of fluid to a solid wall
barrier and second scenario, is to study relation between solid wall and its
collapse time by dam break. The results show sliding distance of a solid wall
is depend on a fluid height, and relation between fluid heght and wall
dimension will be discussed too in this paper.
Keywords: smoothed particle hydrodynamic, dam break, solid-fluid interaction
## 1 Introduction
Dam is built for water supply, water flow or flooding control and electricity
energy storage, but in other hand, dam is one of the most dangerous natural
disaster in many countries including in Indonesia. Because of the age, there
are some evidence of natural disaster by a dam failure and break were reported
for example in Situgintung Tangerang, close to Jakarta was happen in 2009 [1].
The impact of dam break in neighbour area and is huge and many flooding in
remote area, the computer calculation to predict effect of the dam broken is
important including its vast impact[2].
Smoothed Particle Hydrodynamics (SPH), is one of numerical method based on
Lagrangian grid which is applied in astrophysical simulation may be used to
solve the simulation on dam break effect. The development of SPH methods
become alternative methods to solving Navier-Stokes equation, which is main
key in fluid dynamic simulation. As a standard test case in computational
fluid dynamic the dam break simulation were often reported in many cases,
including using SPH methods [5]. The simulation using SPH became emerge
whether using two or three dimension. In other hand, simulation of solid
particle is many reported direct numerical simulation [6]. Interaction between
fluid and solid particle is still open ended to studying. Including simulation
base on SPH method.
In this paper, SPH is developed for supporting solid particles in use for 3D
dam break effect (3D-DBE) simulation. Solid particle have been treated same as
fluid particles with additional calculation for converting gained position
became translation and rotation of solid object in a whole body. With this
capability, the result of 3D-DBE simulation has been varies and interesting.
The goals of this simulation is for analyse fluid and solid particle
interaction by using two different scenario. The first scenario relation
between height of fluid to a solid wall barrier and second scenario, is to
study relation between solid wall and its collapse time by dam break. The
results show sliding distance of a solid wall is depend on a fluid height, and
relation between fluid height and wall dimension will be discussed too.
## 2 Smoothed Particle Hydrodynamics
The governing equation for fluid dynamics is Navier-Stokes equations consist
of conservation of momentum, conservation of mass and conservation of energy.
Assume the system of dam is in isothermal state, then we can neglect the
calculation of fluid’s thermal energy. The momentum equation has been written
as follows:
$\rho\,\frac{\partial\mathbf{u}}{\partial
t}+\rho\,\left(\mathbf{u}\cdot\nabla\,\mathbf{u}\right)=-\nabla
p+\mu\nabla^{2}\mathbf{v}+\rho\mathbf{g}$ (1)
with $\rho$ is fluid density, t is time, v is fluid velocity, p is pressure,
$\mu$ is viscosity, and f is external force. The fluid motion is also governed
by the conservation of mass, which stated by
$\frac{\partial\rho}{\partial t}+\nabla\cdot\left(\rho\mathbf{u}\right)=0$ (2)
Equations (1) and (2) will be solved by numerical approximation of SPH method.
The method simply described by discreet collection of particles or balls for
approximating fluid body. Every particle mathematically influence the
neighbouring particles via kernel function.
$W\left(r,l\right)=\frac{1}{l\sqrt{\pi}}\exp{\left(\frac{-r^{2}}{l^{2}}\right)}$
(3)
Here $W$ is kernel function, $x$ is particle distance and $l$ is smoothing
factor. The kernel function stated in equation (3) now can be use as weighting
function to calculate all physical quantity, gradient and Laplacian attached
to the particle, here we calculate the solution of equation (1) and (2) using
weighted interpolation concept and solved them numerically. The detail of this
method described clearly in [4].
Figure 1: Kernel radius in SPH simulation
Figure 1 illustrated the approximation of fluid body by collection of moveable
particle interacted along kernel radius $r$. The numerical SPH for momentum
and mass conservation int total derivative form respectively are [3]:
$\frac{D\mathbf{u}_{a}}{Dt}=\sum_{b}m_{b}\left(\frac{p_{b}}{{\rho_{b}}^{2}}+\frac{p_{a}}{{\rho_{a}}^{2}}+\Pi_{ab}\right)\nabla_{a}W_{ab}+\mathbf{g}$
(4)
and
$\frac{D{\rho}_{a}}{Dt}=\sum_{b}m_{b}\,\left(\mathbf{u}_{b}-\mathbf{u}_{a}\right)\cdot\nabla_{a}W_{ab}.$
(5)
With index $a$ and $b$ denoted particle index, $\Pi_{ab}$ is numerical
viscosity described in[3]. Here the kernel function $W_{ab}$ and it’s
derivative solved analytically. Equation (4) and (5) can be solved using
direct numerical integration such as Euler method. Another well-known methods
is Runge-Kuta method, Leap frog method and Verlet method.
## 3 Solid - Fluid Interaction
In DBE simulation, one of the important aspect is solid motion exerted by
fluid energy. Here we use solid body dynamic to simulate the motion of fluid.
Once the total force and total torsion exerted by fluid to surface of solid
already known, so do translational and rotational motion of solid can be
calculated.
Figure 2: Force diagram and solid wall motion
Figure 2 show the basic physical scheme of the solid fluid interaction with
strong assumption that no flow across the surface of fluid particle. Some
important aspect such as floor friction and center mass of solid body is also
taking into account.
Figure 3: Solid particle motion schematic
To simplified the continuum aspect of solid body, solid object will be
approximated also by collection rigid particles or balls with certain physical
quantity governed by newton equation of translation and rotation. The center
of mass of such discreet particle are also easy to calculate. Figure 3 showed
the step process of solid particle position update in rotational motion.
## 4 Results and Discussion
The scenario of DBE simulation in three dimension described as follow.
Initially the fluid position in in cubic shape dam for simplification and
without initial velocity of fluid assuming the fluid is at rest. At a certain
distance lied solid wall to block flow of the fluid. Another assumption were
made that there no crack or fracture in solid wall.
Figure 4: Simulation scenario diagram, $h_{2}$ is initial fluid height,
$h_{1}$ and $w_{1}$ is height and width of solid
Figure 4 shown the description of initial arrangement of solid and fluid
particles. As described earlier the arrangement of fluid particle in box shape
dam and initially particle at rest. As the time step moves the gravity force
work on the fluid particles and they move with the same for as if happened in
the dam failure suddenly.
Some parameter involved in 3D simulation in Figure 4 is initial fluid height,
solid wall dimension. Later in this section will be shown some result in 3D
DBE simulation with some variation in simulation parameter.
Figure 5: 3D simulation of DBE a $t_{1}=0s$, $t_{2}=0.228s$, $t_{3}=0.543s$
and $t_{4}=1.068s$, $h_{1}=55m$ and $h_{2}=40m$ and $w_{1}=40m$
Figure 5 show general result of DBE simulation wall with height $55m$ and
fluid dam with $40m$ in height at time $t_{1}=0s$, $t_{2}=0.228s$,
$t_{3}=0.543s$ and $t_{4}=1.068s$. The frame shows the fall process of the
wall as the effect of fluid flow.
Figure 6: Wall height effect on DBE simulation at $t=0.54s$. Left $h_{1}=55m$,
right $h_{1}=60m$
The variation in solid wall height is made to test the stability of moveable
wall under the same energy of fluid flow exerted on the wall surface. It is
done by setting the fluid height at $40m$. Figure 6 shows the effect of wall
height to the fall process under the influence of flow caused by dam break
event. Here the frame plot the position of particles at the same time step
$0.54s$ and the same initial height of fluid $h_{2}=40m$ but in different wall
height. The frames show different response of falling time. Higher wall fall
faster than the lower one at other parameter were keep constant. Here we can
say that the dimension of the wall influence the stability of the wall object
under the same outer forces. If the width of the wall has been increasing the
high wall also become relative stable.
In another scenario the variation of fluid height were tested to the
relatively stable wall object and observe the wall sliding distance.
Figure 7: Fluid height effect to the sliding distance on DBE simulation
Figure 7 shows sliding distance plot versus initial fluid height. The trend of
the data seems to be non linear as the height of fluid increase over 55 m,
this is make sense because of volume of the fluid increase by the power of
three to the height measure and linearly increase the potential energy of
fluid.
## 5 Conclusion
We successfully simulate 3D-DBE simulation using SPH methods, and fluid-solid
interaction capability was applied to study the dam break effect on a solid
wall object, the SPH can be used to simulate a solid particle simulation. Some
simplified assumption that no crack or fracture happen in the solid wall
object thus no separation of solid body as the impact of the fluid momentum.
The simulation results showed that the fluid height and solid wall thickness
and height play a significant effect on solid wall translational and
rotational motion.
For the future work it is also challenging to study the fracture event,
crumbling process or separation of solid body by the effect of momentum from
the fluid flow and made realistic devastation process from dam failure. This
features will increase the computational cost and difficulties in algorithm.
## Acknowledgement
A part of this research was funding by Research and Innovation Program of
Research Group ITB, grant No:243/I.1.C01/PL/2013.
## References
* [1] BBC (2009). Indonesian dam burst toll rises. BBC News Asia-Pacific
* [2] D. Munger (2009). Economic Consequences Methodology for Dam Failure Scenarios. Reclamation Managing in the West,Technical Memorandum Number EC-2009-01 T. U.S. Department of the Interior Bureau of Reclamation Technical Service Center Security, Safety and Law Enforcement
* [3] J.J. Monaghan (1992). Smoothed Particle Hydrodynamic. Annu. Rev. Astron. Astrophys, 30, 543 - 74.
* [4] G.R. Liu and M.B. Liu (2007). Smoothed Particle Hydrodynamics: A Meshfree Particle Method. World Scientific Publishing, Singapore
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|
arxiv-papers
| 2013-09-14T07:35:29 |
2024-09-04T02:49:50.899641
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Suprijadi, Ferry Faizal, Christian Naa, Anggy Trisnawan Putra",
"submitter": "Christian Naa",
"url": "https://arxiv.org/abs/1309.3632"
}
|
1309.3742
|
EUROPEAN ORGANIZATION FOR NUCLEAR RESEARCH (CERN)
CERN-PH-EP-2013-166 LHCb-PAPER-2013-048 November 13, 2013
Measurement of the charge asymmetry in $B^{\pm}\rightarrow\phi K^{\pm}$ and
search for $B^{\pm}\rightarrow\phi\pi^{\pm}$ decays
The LHCb collaboration†††Authors are listed on the following pages.
The $C\\!P$-violating charge asymmetry in $B^{\pm}\rightarrow\phi K^{\pm}$
decays is measured in a sample of $pp$ collisions at 7 TeV centre-of-mass
energy, corresponding to an integrated luminosity of 1.0 fb-1 collected by the
LHCb experiment. The result is $\mathcal{A}_{CP}(B^{\pm}\rightarrow\phi
K^{\pm})=\rm 0.022\pm 0.021\pm 0.009$, where the first uncertainty is
statistical and the second systematic. In addition, a search for the
$B^{\pm}\rightarrow\phi\pi^{\pm}$ decay mode is performed, using the
$B^{\pm}\rightarrow\phi K^{\pm}$ decay rate for normalization. An upper limit
on the branching fraction
$\mathcal{B}(B^{\pm}\rightarrow\phi\pi^{\pm})<1.5\times 10^{-7}$ is set at 90%
confidence level.
Submitted to Phys. Lett. B
© CERN on behalf of the LHCb collaboration, license CC-BY-3.0.
LHCb collaboration
R. Aaij40, B. Adeva36, M. Adinolfi45, C. Adrover6, A. Affolder51, Z.
Ajaltouni5, J. Albrecht9, F. Alessio37, M. Alexander50, S. Ali40, G.
Alkhazov29, P. Alvarez Cartelle36, A.A. Alves Jr24, S. Amato2, S. Amerio21, Y.
Amhis7, L. Anderlini17,f, J. Anderson39, R. Andreassen56, J.E. Andrews57, R.B.
Appleby53, O. Aquines Gutierrez10, F. Archilli18, A. Artamonov34, M. Artuso58,
E. Aslanides6, G. Auriemma24,m, M. Baalouch5, S. Bachmann11, J.J. Back47, A.
Badalov35, C. Baesso59, V. Balagura30, W. Baldini16, R.J. Barlow53, C.
Barschel37, S. Barsuk7, W. Barter46, Th. Bauer40, A. Bay38, J. Beddow50, F.
Bedeschi22, I. Bediaga1, S. Belogurov30, K. Belous34, I. Belyaev30, E. Ben-
Haim8, G. Bencivenni18, S. Benson49, J. Benton45, A. Berezhnoy31, R. Bernet39,
M.-O. Bettler46, M. van Beuzekom40, A. Bien11, S. Bifani44, T. Bird53, A.
Bizzeti17,h, P.M. Bjørnstad53, T. Blake37, F. Blanc38, J. Blouw10, S. Blusk58,
V. Bocci24, A. Bondar33, N. Bondar29, W. Bonivento15, S. Borghi53, A.
Borgia58, T.J.V. Bowcock51, E. Bowen39, C. Bozzi16, T. Brambach9, J. van den
Brand41, J. Bressieux38, D. Brett53, M. Britsch10, T. Britton58, N.H. Brook45,
H. Brown51, A. Bursche39, G. Busetto21,q, J. Buytaert37, S. Cadeddu15, O.
Callot7, M. Calvi20,j, M. Calvo Gomez35,n, A. Camboni35, P. Campana18,37, D.
Campora Perez37, A. Carbone14,c, G. Carboni23,k, R. Cardinale19,i, A.
Cardini15, H. Carranza-Mejia49, L. Carson52, K. Carvalho Akiba2, G. Casse51,
L. Castillo Garcia37, M. Cattaneo37, Ch. Cauet9, R. Cenci57, M. Charles54, Ph.
Charpentier37, P. Chen3,38, S.-F. Cheung54, N. Chiapolini39, M.
Chrzaszcz39,25, K. Ciba37, X. Cid Vidal37, G. Ciezarek52, P.E.L. Clarke49, M.
Clemencic37, H.V. Cliff46, J. Closier37, C. Coca28, V. Coco40, J. Cogan6, E.
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Coombes45, S. Coquereau8, G. Corti37, B. Couturier37, G.A. Cowan49, D.C.
Craik47, M. Cruz Torres59, S. Cunliffe52, R. Currie49, C. D’Ambrosio37, P.
David8, P.N.Y. David40, A. Davis56, I. De Bonis4, K. De Bruyn40, S. De
Capua53, M. De Cian11, J.M. De Miranda1, L. De Paula2, W. De Silva56, P. De
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S. Eisenhardt49, U. Eitschberger9, R. Ekelhof9, L. Eklund50,37, I. El Rifai5,
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Ferguson49, V. Fernandez Albor36, F. Ferreira Rodrigues1, M. Ferro-Luzzi37, S.
Filippov32, M. Fiore16,e, C. Fitzpatrick37, M. Fontana10, F. Fontanelli19,i,
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Furfaro23,k, A. Gallas Torreira36, D. Galli14,c, M. Gandelman2, P. Gandini58,
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Golutvin52,30,37, A. Gomes2, P. Gorbounov30,37, H. Gordon37, M. Grabalosa
Gándara5, R. Graciani Diaz35, L.A. Granado Cardoso37, E. Graugés35, G.
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Grünberg60, B. Gui58, E. Gushchin32, Yu. Guz34,37, T. Gys37, C.
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J. Harrison53, T. Hartmann60, J. He37, T. Head37, V. Heijne40, K. Hennessy51,
P. Henrard5, J.A. Hernando Morata36, E. van Herwijnen37, M. Heß60, A.
Hicheur1, E. Hicks51, D. Hill54, M. Hoballah5, C. Hombach53, W. Hulsbergen40,
P. Hunt54, T. Huse51, N. Hussain54, D. Hutchcroft51, D. Hynds50, V.
Iakovenko43, M. Idzik26, P. Ilten12, R. Jacobsson37, A. Jaeger11, E. Jans40,
P. Jaton38, A. Jawahery57, F. Jing3, M. John54, D. Johnson54, C.R. Jones46, C.
Joram37, B. Jost37, M. Kaballo9, S. Kandybei42, W. Kanso6, M. Karacson37, T.M.
Karbach37, I.R. Kenyon44, T. Ketel41, B. Khanji20, O. Kochebina7, I.
Komarov38, R.F. Koopman41, P. Koppenburg40, M. Korolev31, A. Kozlinskiy40, L.
Kravchuk32, K. Kreplin11, M. Kreps47, G. Krocker11, P. Krokovny33, F. Kruse9,
M. Kucharczyk20,25,37,j, V. Kudryavtsev33, K. Kurek27, T. Kvaratskheliya30,37,
V.N. La Thi38, D. Lacarrere37, G. Lafferty53, A. Lai15, D. Lambert49, R.W.
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Mancinelli6, J. Maratas5, U. Marconi14, P. Marino22,s, R. Märki38, J. Marks11,
G. Martellotti24, A. Martens8, A. Martín Sánchez7, M. Martinelli40, D.
Martinez Santos41,37, D. Martins Tostes2, A. Martynov31, A. Massafferri1, R.
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B. Muryn26, B. Muster38, P. Naik45, T. Nakada38, R. Nandakumar48, I. Nasteva1,
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Palutan18, J. Panman37, A. Papanestis48, M. Pappagallo50, C. Parkes53, C.J.
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Patrignani19,i, C. Pavel-Nicorescu28, A. Pazos Alvarez36, A. Pearce53, A.
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Trigo36, A. Pérez-Calero Yzquierdo35, P. Perret5, M. Perrin-Terrin6, L.
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Zhang58, W.C. Zhang12, Y. Zhang3, A. Zhelezov11, A. Zhokhov30, L. Zhong3, A.
Zvyagin37.
1Centro Brasileiro de Pesquisas Físicas (CBPF), Rio de Janeiro, Brazil
2Universidade Federal do Rio de Janeiro (UFRJ), Rio de Janeiro, Brazil
3Center for High Energy Physics, Tsinghua University, Beijing, China
4LAPP, Université de Savoie, CNRS/IN2P3, Annecy-Le-Vieux, France
5Clermont Université, Université Blaise Pascal, CNRS/IN2P3, LPC, Clermont-
Ferrand, France
6CPPM, Aix-Marseille Université, CNRS/IN2P3, Marseille, France
7LAL, Université Paris-Sud, CNRS/IN2P3, Orsay, France
8LPNHE, Université Pierre et Marie Curie, Université Paris Diderot,
CNRS/IN2P3, Paris, France
9Fakultät Physik, Technische Universität Dortmund, Dortmund, Germany
10Max-Planck-Institut für Kernphysik (MPIK), Heidelberg, Germany
11Physikalisches Institut, Ruprecht-Karls-Universität Heidelberg, Heidelberg,
Germany
12School of Physics, University College Dublin, Dublin, Ireland
13Sezione INFN di Bari, Bari, Italy
14Sezione INFN di Bologna, Bologna, Italy
15Sezione INFN di Cagliari, Cagliari, Italy
16Sezione INFN di Ferrara, Ferrara, Italy
17Sezione INFN di Firenze, Firenze, Italy
18Laboratori Nazionali dell’INFN di Frascati, Frascati, Italy
19Sezione INFN di Genova, Genova, Italy
20Sezione INFN di Milano Bicocca, Milano, Italy
21Sezione INFN di Padova, Padova, Italy
22Sezione INFN di Pisa, Pisa, Italy
23Sezione INFN di Roma Tor Vergata, Roma, Italy
24Sezione INFN di Roma La Sapienza, Roma, Italy
25Henryk Niewodniczanski Institute of Nuclear Physics Polish Academy of
Sciences, Kraków, Poland
26AGH - University of Science and Technology, Faculty of Physics and Applied
Computer Science, Kraków, Poland
27National Center for Nuclear Research (NCBJ), Warsaw, Poland
28Horia Hulubei National Institute of Physics and Nuclear Engineering,
Bucharest-Magurele, Romania
29Petersburg Nuclear Physics Institute (PNPI), Gatchina, Russia
30Institute of Theoretical and Experimental Physics (ITEP), Moscow, Russia
31Institute of Nuclear Physics, Moscow State University (SINP MSU), Moscow,
Russia
32Institute for Nuclear Research of the Russian Academy of Sciences (INR RAN),
Moscow, Russia
33Budker Institute of Nuclear Physics (SB RAS) and Novosibirsk State
University, Novosibirsk, Russia
34Institute for High Energy Physics (IHEP), Protvino, Russia
35Universitat de Barcelona, Barcelona, Spain
36Universidad de Santiago de Compostela, Santiago de Compostela, Spain
37European Organization for Nuclear Research (CERN), Geneva, Switzerland
38Ecole Polytechnique Fédérale de Lausanne (EPFL), Lausanne, Switzerland
39Physik-Institut, Universität Zürich, Zürich, Switzerland
40Nikhef National Institute for Subatomic Physics, Amsterdam, The Netherlands
41Nikhef National Institute for Subatomic Physics and VU University Amsterdam,
Amsterdam, The Netherlands
42NSC Kharkiv Institute of Physics and Technology (NSC KIPT), Kharkiv, Ukraine
43Institute for Nuclear Research of the National Academy of Sciences (KINR),
Kyiv, Ukraine
44University of Birmingham, Birmingham, United Kingdom
45H.H. Wills Physics Laboratory, University of Bristol, Bristol, United
Kingdom
46Cavendish Laboratory, University of Cambridge, Cambridge, United Kingdom
47Department of Physics, University of Warwick, Coventry, United Kingdom
48STFC Rutherford Appleton Laboratory, Didcot, United Kingdom
49School of Physics and Astronomy, University of Edinburgh, Edinburgh, United
Kingdom
50School of Physics and Astronomy, University of Glasgow, Glasgow, United
Kingdom
51Oliver Lodge Laboratory, University of Liverpool, Liverpool, United Kingdom
52Imperial College London, London, United Kingdom
53School of Physics and Astronomy, University of Manchester, Manchester,
United Kingdom
54Department of Physics, University of Oxford, Oxford, United Kingdom
55Massachusetts Institute of Technology, Cambridge, MA, United States
56University of Cincinnati, Cincinnati, OH, United States
57University of Maryland, College Park, MD, United States
58Syracuse University, Syracuse, NY, United States
59Pontifícia Universidade Católica do Rio de Janeiro (PUC-Rio), Rio de
Janeiro, Brazil, associated to 2
60Institut für Physik, Universität Rostock, Rostock, Germany, associated to 11
61Celal Bayar University, Manisa, Turkey, associated to 37
aP.N. Lebedev Physical Institute, Russian Academy of Science (LPI RAS),
Moscow, Russia
bUniversità di Bari, Bari, Italy
cUniversità di Bologna, Bologna, Italy
dUniversità di Cagliari, Cagliari, Italy
eUniversità di Ferrara, Ferrara, Italy
fUniversità di Firenze, Firenze, Italy
gUniversità di Urbino, Urbino, Italy
hUniversità di Modena e Reggio Emilia, Modena, Italy
iUniversità di Genova, Genova, Italy
jUniversità di Milano Bicocca, Milano, Italy
kUniversità di Roma Tor Vergata, Roma, Italy
lUniversità di Roma La Sapienza, Roma, Italy
mUniversità della Basilicata, Potenza, Italy
nLIFAELS, La Salle, Universitat Ramon Llull, Barcelona, Spain
oHanoi University of Science, Hanoi, Viet Nam
pInstitute of Physics and Technology, Moscow, Russia
qUniversità di Padova, Padova, Italy
rUniversità di Pisa, Pisa, Italy
sScuola Normale Superiore, Pisa, Italy
## 1 Introduction
The weak-interaction $B^{\pm}\\!\rightarrow\phi K^{\pm}$ decay is governed by
the $b\\!\rightarrow ss\overline{}s$ transition. In the Standard Model (SM),
it can only occur through loop diagrams (see Fig. 1),
Figure 1: Lowest-order Feynman diagrams of the Standard Model for the decays
$B^{+}\\!\rightarrow\phi K^{+}$ (top) and $B^{+}\\!\rightarrow\phi\pi^{+}$
(bottom). The diagrams with an external $\phi$ meson are OZI suppressed.
leading to a branching fraction of order $10^{-5}$ [1]. Because the dominant
amplitudes have similar weak phases, the $C\\!P$-violating charge asymmetry,
defined as
$\mathcal{A}_{C\\!P}(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})\equiv\frac{{\cal B}(\mbox{$B^{-}\\!\rightarrow\phi K^{-}$})-{\cal
B}(\mbox{$B^{+}\\!\rightarrow\phi K^{+}$})}{{\cal
B}(\mbox{$B^{-}\\!\rightarrow\phi K^{-}$})+{\cal
B}(\mbox{$B^{+}\\!\rightarrow\phi K^{+}$})}\,,$ (1)
is predicted to be small in the SM, typically 1–2% with uncertainties of a few
percent [2, 3]. A significantly larger value would signal interference with an
amplitude not described in the SM. The current experimental world average is
$\mathcal{A}_{C\\!P}(\mbox{\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$}})=0.10\pm 0.04$ [1], dominated by a recent measurement from the
BaBar collaboration [4]. Large $CP$ violation effects have been seen in some
regions of the $B^{\pm}\\!\rightarrow K^{+}K^{-}K^{\pm}$ phase space, but not
around the $\phi$ resonance [5].
The $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ decay is another flavour-changing
neutral current process, driven by the $b\\!\rightarrow ds\overline{}s$ quark-
level transition (see Fig. 1). The high suppression, due to the tiny product
of the Cabibbo-Kobayashi-Maskawa matrix elements [6, 7] and to the Okubo-
Zweig-Iizuka (OZI) rule [8, 9, 10] associated with the creation of the
colourless $s\overline{}s$ pair forming the $\phi$ meson, makes this rare loop
decay a sensitive probe of the SM. Indeed, even a small non-SM amplitude, e.g.
from R-parity violating supersymmetry [11], may dominate over the SM
contribution.
The current SM prediction for the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$
branching fraction suffers from uncertainties originating from the naïve
factorization approach, radiative corrections, calculation of the long-
distance contribution (e.g. $B\\!\rightarrow KK^{*}$ rescattering), and
$\omega-\phi$ mixing [12]. The latter is the main source of uncertainty. The
physical $\omega$ and $\phi$ meson states do not coincide exactly with the
ideal $(|u\overline{}u\rangle+|d\overline{}d\rangle)/\sqrt{2}$ and
$|s\overline{}s\rangle$ states, respectively. They appear to be mixtures of
these two states characterized by a small mixing angle $\delta_{V}$ [13, 14],
which depends on the magnitude of $\mathrm{SU}(3)$ symmetry breaking and can
be determined in the framework of chiral perturbation theory. However, more
sophisticated treatments based on the full $\rho^{0}-\omega-\phi$ mixing
scheme suggest that $\delta_{V}$ is mass dependent, i.e. takes different
values at the $\omega$ and $\phi$ masses [15, 16]. In the QCD factorization
approach, the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ branching fraction is
predicted to be in the range $(5-10)\times 10^{-9}$ [3] if $\omega-\phi$
mixing is neglected, but can be enhanced up to $0.6\times 10^{-7}$ [17, 12]
depending on the value of $\delta_{V}$. However, the effect of $\omega-\phi$
mixing has not been observed in a recent search for
$B^{0}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}\phi$ [18].
Values of the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ branching fraction in
excess of $10^{-7}$ would be indicative of non-SM physics.
The $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ decay mode has not been observed yet.
Currently, the most stringent experimental limit is ${\cal
B}(\mbox{$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$})<2.4\times 10^{-7}$ at 90%
confidence level (CL), obtained by the BaBar collaboration [19].
This Letter presents a measurement of the $B^{\pm}\\!\rightarrow\phi K^{\pm}$
charge asymmetry and a search for the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$
decay mode with the LHCb detector. The results are based on a data sample
collected during the 2011 $pp$ run of the Large Hadron Collider at a centre-
of-mass energy of 7$\mathrm{\,Te\kern-1.00006ptV}$, corresponding to an
integrated luminosity of 1.0$\mbox{\,fb}^{-1}$. The $\phi$ meson is
reconstructed in the $K^{+}K^{-}$ final state. We define the $\phi$ signal as
any peaking component in the $K^{+}K^{-}$ mass spectrum consistent with the
known parameters of the $\phi$ resonance, without attempting a full amplitude
analysis of the three-body $K^{+}K^{-}K^{\pm}$ and $K^{+}K^{-}\pi^{\pm}$ final
states. In order to suppress several systematic effects, the primary
observables measured in this analysis are the difference of $C\\!P$-violating
charge asymmetries
$\Delta\mathcal{A}_{C\\!P}\equiv\mathcal{A}_{C\\!P}(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})-\mathcal{A}_{C\\!P}(\mbox{$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}K^{\pm}$})\,,$ (2)
and the branching fraction ratio ${\cal
B}(\mbox{$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$})/{\cal
B}(\mbox{$B^{\pm}\\!\rightarrow\phi K^{\pm}$})$, which are then converted to
results on $\mathcal{A}_{C\\!P}(\mbox{$B^{\pm}\\!\rightarrow\phi K^{\pm}$})$
and ${\cal B}(\mbox{$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$})$ using the best
known values of
$\mathcal{A}_{C\\!P}(\mbox{$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}K^{\pm}$})$ [1, 20] and ${\cal B}(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})$ [1]. The choice of
$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$ as
reference channel and other features of the analysis follow the approach
adopted in inclusive studies of $B^{\pm}\\!\rightarrow K^{+}K^{-}K^{\pm}$
decays with the same data set [5].
The two measurements are performed in a common analysis, i.e. they are based
on identical event selections and data descriptions whenever possible. The
observables are obtained from two-dimensional maximum likelihood fits to the
unbinned $B^{\pm}$ and $\phi$ mass distributions of the reconstructed
candidates, using parametric shapes with minimal dependence on simulation. The
results of these fits were not examined until the entire analysis procedure
was finalized.
## 2 Detector and data set
The LHCb detector [21] is a single-arm forward spectrometer covering the
pseudorapidity range $2<\eta<5$, designed for the study of particles
containing $b$ or $c$ quarks. The detector includes a high-precision tracking
system consisting of a silicon-strip vertex detector surrounding the $pp$
interaction region, a large-area silicon-strip detector located upstream of a
dipole magnet with a bending power of about $4{\rm\,Tm}$, and three stations
of silicon-strip detectors and straw drift tubes placed downstream. The
combined tracking system provides a momentum measurement with a relative
uncertainty that varies from 0.4% at 5${\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$
to 0.6% at 100${\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$, and an impact parameter
(IP) resolution of 20$\,\upmu\rm m$ for tracks with high transverse momentum
($p_{\rm T}$). Charged hadrons are identified using two ring-imaging Cherenkov
detectors [22]. Photon, electron and hadron candidates are identified by a
calorimeter system consisting of scintillating-pad and preshower detectors, an
electromagnetic calorimeter and a hadronic calorimeter. Muons are identified
by a system composed of alternating layers of iron and multiwire proportional
chambers. The direction of the magnetic field of the spectrometer dipole
magnet is reversed regularly.
The trigger [23] consists of a hardware stage, based on information from the
calorimeter and muon systems, followed by a software stage, which applies a
full event reconstruction. The $B^{\pm}$ candidate decays considered in this
analysis must belong to one of two exclusive categories of events, called TOS
(triggered on signal) or TIS (triggered independently of signal). A TOS event
is triggered at the hardware stage by one of the candidate’s final-state
particles being compatible with a transverse energy deposit greater than
3.5$\mathrm{\,Ge\kern-1.00006ptV}$ in the hadron calorimeter. A TIS event does
not satisfy the TOS definition and is triggered at the hardware stage by
activity in the rest of the event. All candidates must pass a software trigger
requiring a two-, three- or four-track secondary vertex with a large scalar
sum of the transverse momentum of the tracks and a significant displacement
from the primary $pp$ interaction vertices (PVs). At least one track should
have $\mbox{$p_{\rm T}$}>1.7{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$ and
$\chi^{2}_{\rm IP}$ with respect to any PV greater than 16, where
$\chi^{2}_{\rm IP}$ is defined as the difference in $\chi^{2}$ of a given PV
reconstructed with and without the considered track. A multivariate algorithm
[24] is used for the identification of secondary vertices consistent with the
decay of a $b$ hadron.
In the simulation, $pp$ collisions are generated using Pythia 6.4 [25] with a
specific LHCb configuration [26]. Decays of hadronic particles are described
by EvtGen [27], in which final state radiation is generated using Photos [28].
The interaction of the generated particles with the detector and its response
are implemented using the Geant4 toolkit [29, *Agostinelli:2002hh] as
described in Ref. [31].
## 3 Event selection and efficiency
The selections of $B^{\pm}\\!\rightarrow\phi K^{\pm}$ and
$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ candidates are identical, except for the
particle identification (PID) requirement on the charged hadron combined with
the $\phi$ candidate, which is referred to as the bachelor hadron $h^{\pm}$
($h^{\pm}=K^{\pm}$ or $\pi^{\pm}$). The other requirements are chosen to
minimize the relative statistical uncertainty on the
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ signal yield.
Only good quality tracks with $\chi^{2}_{\rm IP}>25$ and $\mbox{$p_{\rm
T}$}>0.25{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$ are used in the reconstruction.
The $\phi$ meson candidates are reconstructed from two oppositely-charged
tracks identified as kaons with the PID requirement ${\rm DLL}_{K\pi}>2$,
where ${\rm DLL}_{K\pi}$ is the difference in log-likelihood between the kaon
and pion hypotheses, as determined with the ring-imaging Cherenkov detectors
in control samples of known particle composition [22]. The $\phi$ candidates
are required to have $\mbox{$p_{\rm
T}$}>2{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$, a total momentum, $p$, larger
than 10${\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$ and an invariant mass, $m_{KK}$,
in the range 1.00–1.05${\mathrm{\,Ge\kern-1.00006ptV\\!/}c^{2}}$. Bachelor
hadrons, reconstructed either as pions if ${\rm DLL}_{K\pi}<-1$ or kaons
otherwise, are required to have $p>10{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$ and
$\mbox{$p_{\rm T}$}>2.5{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$, and are combined
with $\phi$ candidates to form $B^{\pm}\\!\rightarrow\phi h^{\pm}$ candidates.
These $B^{\pm}$ candidates are required to have $\mbox{$p_{\rm
T}$}>2{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$, a three-track vertex $\chi^{2}$
per degree of freedom less than 9, and an invariant mass $m_{KKh}$ in the
range 5.0–5.5${\mathrm{\,Ge\kern-1.00006ptV\\!/}c^{2}}$. Furthermore
$\cos\theta_{\rm p}$ is required to be greater than 0.9999, where $\theta_{\rm
p}$ is the angle between the $B^{\pm}$ momentum vector and the vector joining
the $B^{\pm}$ production vertex to the $B^{\pm}$ decay vertex. The production
vertex is chosen as the PV for which the $B^{\pm}$ has the smallest
$\chi^{2}_{\rm IP}$.
Multiple candidates, occurring in 0.2% of the events, are removed by keeping
the candidate with the smallest $B^{\pm}$ vertex $\chi^{2}$. The final data
sample consists of 6251 $B^{\pm}\\!\rightarrow\phi K^{\pm}$ candidates and
2169 $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ candidates.
The PID performance is determined from a large and high-purity sample of pions
from prompt $D^{*+}\\!\rightarrow D^{0}(K^{-}\pi^{+})\pi^{+}$ and
$D^{*-}\\!\rightarrow\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}(K^{+}\pi^{-})\pi^{-}$ decays, as a
function of $p$ and $\eta$. After reweighting this calibration sample to the
same momentum and pseudorapidity distributions as for the bachelor pion in
simulated $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ decays, the efficiency of the
PID requirement ${\rm DLL}_{K\pi}<-1$ for the bachelor pion is measured to be
$\rm 0.846\pm 0.011\,(stat)\pm 0.020\,(syst)$, with a 5% kaon
misidentification probability. All other efficiencies, which are slightly
different for $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ and
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ decays due to their kinematic properties,
are determined from simulation. The efficiency ratio
$\frac{\epsilon(\mbox{$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$})}{\epsilon(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})}=\rm 0.762\pm 0.031\,(stat)\pm 0.036\,(syst)$ (3)
is obtained, where the numerator is the total efficiency for a
$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ decay to be selected as a
$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ candidate and the denominator is the
total efficiency for a $B^{\pm}\\!\rightarrow\phi K^{\pm}$ decay to be
selected either as a $B^{\pm}\\!\rightarrow\phi K^{\pm}$ candidate or as a
$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ candidate. The statistical uncertainty
arises from the size of the calibration and simulation samples, while the
systematic uncertainty is the quadratic sum of contributions from the PID
($\pm 0.018$), the trigger ($\pm 0.008$), and other offline kinematic
selection requirements ($\pm 0.030)$.
## 4 Fit description
The observables of interest, namely the asymmetry between the
$B^{-}\\!\rightarrow\phi K^{-}$ and $B^{+}\\!\rightarrow\phi K^{+}$ yields and
the ratio between the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ and
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ yields, are each determined from a two-
dimensional unbinned extended maximum likelihood fit based on probability
density functions (PDFs) of the $m_{KKh}$ and $m_{KK}$ masses. In each case,
independent subsamples of events, each with either $B^{\pm}\\!\rightarrow\phi
K^{\pm}$ candidates or $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ candidates, are
fitted simultaneously. For each subsample, the likelihood is written as
${\cal
L}={\textstyle\exp(-\sum_{j}N_{j})}\,\prod_{i}^{N}{\left(\textstyle\sum_{j}N_{j}P_{j}^{i}\right)}\,,$
(4)
where $N_{j}$ is the yield of fit component $j$, $P_{j}^{i}$ is the
probability of event $i$ for component $j$, and the index $i$ runs over the
$N$ events in the subsample. Except for the misidentified components described
further below, the probabilities $P_{j}^{i}$ are given by the product of PDFs
for the two $K^{+}K^{-}h^{\pm}$ and $K^{+}K^{-}$ invariant masses, evaluated
at the values $m_{KKh}^{i}$ and $m_{KK}^{i}$ of event $i$:
$P_{j}^{i}=P_{j}^{KKh}(m_{KKh}^{i})P_{j}^{KK}(m_{KK}^{i})\,.$ (5)
This assumes that the two mass variables are independent, as supported by data
and simulation studies. The correlation between $m_{KKh}$ and $m_{KK}$ is
found to be less than $4\%$.
The description of the $m_{KKh}$ distributions involves a combination of three
contributions: a signal peaking at the $B^{\pm}$ mass, a broad low-mass
background with an end-point near
$5150{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ due to partially-reconstructed
$b$-hadron decays such as $B^{0}\\!\rightarrow\phi K^{*0}$, and a linear
background from random combinations. The peaking signal is modelled with a
Crystal Ball function [32] modified such that both the upper and lower tails
are power laws. The mean and the width $\sigma_{B}$ of the Crystal Ball
function are free in the fit, while the tail parameters are determined from
simulation. The partially-reconstructed background is described with an ARGUS
function [33] convoluted with a Gaussian resolution function of the same
$\sigma_{B}$ as the $B^{\pm}$ signal. The $m_{KK}$ distribution is described
with two contributions: a peaking term centred on the $\phi$ mass, described
with a relativistic Breit-Wigner function convoluted with a Gaussian
resolution function of free width, and a linear term originating from
nonresonant, S-wave, or random combinations of two kaons. The above three
$m_{KKh}$ contributions and two $m_{KK}$ contributions lead to six components
for each subsample: the $B^{\pm}\\!\rightarrow\phi h^{\pm}$ signal, the
nonresonant $B^{\pm}\\!\rightarrow K^{+}K^{-}h^{\pm}$ background, the
partially-reconstructed $b$-hadron backgrounds with or without a true $\phi$
meson (for example $B\\!\rightarrow\phi h^{\pm}\pi$ or $B\\!\rightarrow
K^{+}K^{-}h^{\pm}\pi$), and the combinatorial backgrounds with or without a
true $\phi$ meson. The nonresonant $B^{\pm}\\!\rightarrow K^{+}K^{-}h^{\pm}$
components include $b\rightarrow c$ decays, which are found to be negligible
from simulation studies.
In addition, we consider components for the misidentified
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ and $B^{\pm}\\!\rightarrow
K^{+}K^{-}K^{\pm}$ decays in the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ sample,
while misidentified $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ and
$B^{\pm}\\!\rightarrow K^{+}K^{-}\pi^{\pm}$ decays in the
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ sample are negligible, and therefore
ignored. For these two additional components, the $m_{KK\pi}$ PDF is
conditional to the observable $\delta m=m_{KKK}-m_{KK\pi}$, which is the mass
difference under the two bachelor hadron mass hypotheses. The probabilities
are written as
$P_{j}^{i}=P_{\rm misID}^{KK\pi}(m_{KK\pi}^{i}|\delta
m^{i})P_{j}^{KK}(m_{KK}^{i})\,,$ (6)
where $P_{j}^{KK}$ is the $m_{KK}$ PDF described above (representing either
$\phi$ signal or background) and
$\left.P_{\rm misID}^{KK\pi}(m^{i}_{KK\pi}|\delta m^{i})=P_{\phi
K}^{KKK}(m^{i}_{KK\pi}+\delta
m^{i})\right|_{\sigma_{B}\rightarrow\rho\sigma_{B}}\,.$ (7)
Here $P_{\phi K}^{KKK}$ is the $m_{KKK}$ PDF of the $B^{\pm}\\!\rightarrow\phi
K^{\pm}$ signal, but with an increased $B^{\pm}$ mass resolution to account
for the effects of the typically higher momentum of misidentified bachelor
kaons. The parameter $\sigma_{B}$ is multiplied here by the central value of a
factor $\rho=1.26\pm 0.10$, determined from data as the ratio of the measured
$m_{KKK}$ resolutions of the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ signal in the
regions $-7<{\rm DLL}_{K\pi}<-1$ and ${\rm DLL}_{K\pi}>-1$. The expression in
Eq. 7 is equivalent to $\left.P_{\phi
K}^{KKK}(m^{i}_{KKK})\right|_{\sigma_{B}\rightarrow\rho\sigma_{B}}$, which
means that the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ misidentified component in
the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ sample would have a
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ signal distribution if the correct mass
was assigned to the bachelor kaon. The advantage of introducing the $\delta m$
observable is to connect the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ shapes in the
$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ and $B^{\pm}\\!\rightarrow\phi K^{\pm}$
samples, thereby constraining the misidentified $B^{\pm}\\!\rightarrow\phi
K^{\pm}$ component in the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ sample using
the large signal in the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ sample. This
procedure allows to describe the misidentified component with the same
parametric shape as the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ signal, and
reduces the statistical uncertainty on the
$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ yield by a factor of two. However, this
introduces a bias because the $\delta m$ distribution, which is not accounted
for in the likelihood, is not the same for all components [34]. To reduce this
bias, the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ sample is divided into four
bins of $\delta m$, each with its own eight components. This procedure reduces
the bias on the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ signal yield to a
negligible level.
Other fit parameters that are common to the different subsamples are the
$m_{KKh}$ end-point of the partially-reconstructed backgrounds, the peaking
$m_{KK}$ PDF parameters for all components containing a $\phi$ meson, and the
$m_{KK}$ slope of the nonresonant $B^{\pm}\\!\rightarrow K^{+}K^{-}h^{\pm}$
components. Finally, the ratio of the yield of the misidentified nonresonant
$B^{\pm}\\!\rightarrow K^{+}K^{-}K^{\pm}$ background to the yield of
misidentified $B^{\pm}\\!\rightarrow\phi K^{\pm}$ background in the
$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ sample is constrained to the yield ratio
of the corresponding correctly-identified components in the
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ sample.
The fit procedure is validated on simulated data containing the expected
proportion of signal and background events.
These studies, which take into account the different $\delta m$ distributions
and the possible correlation between the fit observables, demonstrate the
stability of the fit and show that the fit results are unbiased.
## 5 Measurement of the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ charge asymmetry
The charge asymmetry of the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ signal is
determined from a fit to the $B^{-}\\!\rightarrow\phi K^{-}$ and
$B^{+}\\!\rightarrow\phi K^{+}$ candidates in the ${\rm DLL}_{K\pi}>-1$
region. These two samples are each divided into two subsamples according to
whether the events were TOS or TIS at the hardware trigger stage. In the fit,
each of the six components has therefore four yields. For the signal
component, they are expressed as $N^{\pm}_{\rm TOS}=N_{\rm
TOS}(1\mp\mathcal{A}_{\rm raw,TOS})/2$ and $N^{\pm}_{\rm TIS}=N_{\rm
TIS}(1\mp\mathcal{A}_{\rm raw,TIS})/2$, where $N_{k}$ is the total yield and
$\mathcal{A}_{{\rm raw},k}$ is the raw yield asymmetry in subsample $k$
($k=\rm TOS,TIS$). The fit has a total of 34 free parameters: 10 mass shape
parameters, 12 yields and 12 raw asymmetries.
\begin{overpic}[width=218.51521pt]{BMass_Asymm_neg.pdf}\put(85.0,73.0){\scriptsize(a)}
\end{overpic}\begin{overpic}[width=218.51521pt]{PhiMass_Asymm_neg.pdf}
\put(85.0,73.0){\scriptsize(b)}
\end{overpic}\begin{overpic}[width=218.51521pt]{BMass_Asymm_pos.pdf}\put(85.0,73.0){\scriptsize(c)}
\end{overpic}\begin{overpic}[width=218.51521pt]{PhiMass_Asymm_pos.pdf}
\put(85.0,73.0){\scriptsize(d)} \end{overpic}
Figure 2: Distributions of the (a) $K^{+}K^{-}K^{-}$ and (b) $K^{+}K^{-}$
masses of the selected $B^{-}\\!\rightarrow\phi K^{-}$ candidates, as well as
of the (c) $K^{+}K^{-}K^{+}$ and (d) $K^{+}K^{-}$ masses of the selected
$B^{+}\\!\rightarrow\phi K^{+}$ candidates. The solid blue curves represent
the result of the simultaneous fit described in the text, with the following
components: $B^{\pm}\\!\rightarrow\phi K^{\pm}$ signal (dotted red),
nonresonant $B^{\pm}\\!\rightarrow K^{+}K^{-}K^{\pm}$ background (dashed red),
partially-reconstructed $b$-hadron background with (dotted blue) or without
(dashed blue) a true $\phi$ meson, and combinatorial background with (dotted
green) or without (dashed green) a true $\phi$ meson. Some of the components
are barely visible because the corresponding yields are small. Normalized
residuals are displayed below each histogram.
Figure 2 shows the projections of the fitting function superimposed on the
$m_{KKK}$ and $m_{KK}$ distributions, shown separately for $B^{-}$ and $B^{+}$
candidates, but where TOS and TIS events are summed. The $m_{KKK}$ resolution
measured from the fit is $\sigma_{B}=20.4\pm
0.3{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$. The fitted raw asymmetries for
the signal are shown in the first line of Table 1. They are statistically
uncorrelated.
Table 1: Raw charge asymmetries for the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ and $B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$ decays, their difference $\Delta\mathcal{A}_{C\\!P}$, and the fraction of $B^{\pm}\\!\rightarrow\phi K^{\pm}$ signal events in each trigger subsample for $k=\rm TOS,TIS$. All uncertainties are statistical only. | TOS subsample | TIS subsample
---|---|---
$\mathcal{A}_{{\rm raw},k}(\mbox{$B^{\pm}\\!\rightarrow\phi K^{\pm}$})$ | $+0.027\pm 0.026$ | $-0.053\pm 0.035$
$\mathcal{A}_{{\rm raw},k}(\mbox{$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$})$ | $-0.024\pm 0.008$ | $-0.008\pm 0.005$
$\Delta\mathcal{A}_{C\\!P}$ | $+0.052\pm 0.027$ | $-0.045\pm 0.035$
$N_{k}/(N_{\rm TOS}+N_{\rm TIS})$ | 66% | 34%
Weighted $\Delta\mathcal{A}_{C\\!P}$ average | $+0.019\pm 0.021$
Each raw charge asymmetry is related to the $C\\!P$ asymmetry through
$\mathcal{A}_{{\rm raw},k}(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})=\mathcal{A}_{C\\!P}(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})+{\cal A}_{{\rm D},k}(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})+{\cal A}_{\rm P}\,,$ (8)
where ${\cal A}_{{\rm D},k}(\mbox{$B^{\pm}\\!\rightarrow\phi K^{\pm}$})$ is
the detection charge asymmetry for the bachelor $K^{\pm}$ and ${\cal A}_{\rm
P}$ is the production asymmetry of $B^{\pm}$ mesons. Equation 8 and the
corresponding equation for the
$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$
reference channel hold because all involved asymmetries are small. Under the
assumption that the detection asymmetry is the same for
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ and
$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$,
which is correct in the limit where the bachelor $K^{\pm}$ has the same
kinematic properties, the difference in charge asymmetries defined in Eq. 2
can be written as
$\Delta\mathcal{A}_{C\\!P}=\mathcal{A}_{{\rm
raw},k}(\mbox{$B^{\pm}\\!\rightarrow\phi K^{\pm}$})-\mathcal{A}_{{\rm
raw},k}(\mbox{$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}K^{\pm}$})$ (9)
and should not depend on the trigger category $k$. The raw charge asymmetries
of $B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$
decays have been measured in a previous analysis [5]; they are subtracted from
the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ raw asymmetries to obtain two
independent measurements of $\Delta\mathcal{A}_{C\\!P}$. Since the two results
agree within about two statistical standard deviations, the results are
combined. The final $\Delta\mathcal{A}_{C\\!P}$ result is computed as a
weighted average, with weights equal to the fractions $N_{k}/(N_{\rm
TOS}+N_{\rm TIS})$ of signal events in the two trigger subsamples. All inputs
to the calculation are reported in Table 1. The separation between TIS and TOS
events is needed because the detection asymmetry ${\cal A}_{{\rm D},k}$
depends on the trigger category $k$ and the fraction of events in the two
categories differs between the signal and reference channels.
Table 2: Systematic uncertainties on the measurement of $\Delta\mathcal{A}_{C\\!P}$. Source | Uncertainty
---|---
Mass shape modelling | 0.003
Possible S-wave contribution | 0.002
Trigger | 0.004
Bachelor kaon kinematic properties | 0.005
Geometric acceptance | 0.002
Quadratic sum | 0.007
Several systematic uncertainties are considered on the weighted
$\Delta\mathcal{A}_{C\\!P}$ average, as summarized in Table 2. The
contribution due to the mass shape modelling is obtained by repeating the fit
(and the calculation of Table 1) with the fixed parameter values of the
Crystal Ball and ARGUS functions changed within their uncertainties, as
determined from simulation and $B^{\pm}\\!\rightarrow\phi K^{\pm}$ data,
respectively, or with an exponential (rather than linear) combinatorial
background model. Possible residual effects from S-wave contributions not
fully accounted for by the linear component are investigated by comparing the
observed angular distribution of the $B^{\pm}\\!\rightarrow\phi K^{\pm}$
signal with the expectation for a peaking structure in the $K^{+}K^{-}$ mass
due to a single P-wave state. Other P-wave components are neglected. If these
S-wave contributions corresponded to an additional component included in the
signal without charge asymmetry, a bias would appear on
$\Delta\mathcal{A}_{C\\!P}$, which is taken as a systematic uncertainty.
The charge asymmetry in the trigger efficiency for kaons of the TOS subsample
does not completely cancel in $\Delta\mathcal{A}_{C\\!P}$, because of the
different number of kaons in the two decay modes considered. The difference
between the values of $\mathcal{A}_{\rm
raw,TOS}(\mbox{$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}K^{\pm}$})$ computed with and without a charge-dependent correction for
the kaon efficiency determined from calibration data is propagated as a
systematic uncertainty on $\Delta\mathcal{A}_{C\\!P}$. Such an effect is
absent for the TIS subsample. Another small contribution, due to the TOS
events that would still be accepted by the hardware trigger level without
considering the particles from the $B^{+}$ candidate decay, has been included
in the trigger systematic uncertainty. Due to differences in the kinematic
selections of the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ and
$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$
decay modes, the assumption of Eq. 9 cannot be exact, and a further systematic
uncertainty is assigned. The fit of the raw charge asymmetries of
$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$ is
repeated with the same kinematic selection on the bachelor kaon as for
$B^{\pm}\\!\rightarrow\phi K^{\pm}$, i.e.
$p>10{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$ and $\mbox{$p_{\rm
T}$}>2.5{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$, and after reweighting its
momentum distribution to that observed in the $B^{\pm}\\!\rightarrow\phi
K^{\pm}$ decays. The resulting effect on $\Delta\mathcal{A}_{C\\!P}$ is taken
as a systematic uncertainty. Finally, we repeat the $B^{\pm}\\!\rightarrow\phi
K^{\pm}$ analysis after requiring the bachelor kaon momentum to point in a
fiducial solid angle avoiding detector edge effects, and assign the observed
change in $\Delta\mathcal{A}_{C\\!P}$ as a systematic uncertainty due to the
geometrical acceptance.
The final measurement is
$\Delta\mathcal{A}_{C\\!P}=\rm 0.019\pm 0.021(stat)\pm 0.007(syst)\,.$ (10)
A recent update of the
$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$
charge asymmetry measurement by the D0 collaboration [20] has not been
included yet in the average of the Particle Data Group (PDG) [1]. Replacing
the previous D0 result with the new one yields the world average
$\mathcal{A}_{C\\!P}(\mbox{$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}K^{\pm}$})=0.003\pm 0.006$, where the uncertainty is scaled by a factor
1.8 according to the PDG averaging rules. Using this average, we obtain
$\mathcal{A}_{C\\!P}(\mbox{$B^{\pm}\\!\rightarrow\phi K^{\pm}$})=\rm 0.022\pm
0.021(stat)\pm 0.009(syst)\,,$ (11)
where the uncertainty on the
$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$
charge asymmetry is incorporated in the systematic uncertainty.
## 6 Search for $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ decays
\begin{overpic}[width=218.51521pt]{BMassPi_PHIPI_Linear_Pad_cut.pdf}\put(85.0,51.0){\scriptsize(a)}
\end{overpic}\begin{overpic}[width=218.51521pt]{PhiMassPi_PHIPI_Linear_Pad_cut.pdf}
\put(85.0,51.0){\scriptsize(b)}
\end{overpic}\begin{overpic}[width=218.51521pt]{BMassPi_PHIPI_Log.pdf}\put(85.0,73.0){\scriptsize(c)}
\end{overpic}\begin{overpic}[width=218.51521pt]{PhiMassPi_PHIPI_Log.pdf}
\put(85.0,73.0){\scriptsize(d)} \end{overpic}
Figure 3: Distributions of the (a, c) $K^{+}K^{-}\pi^{\pm}$ and (b, d)
$K^{+}K^{-}$ masses of the selected $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$
candidates, shown both with linear and logarithmic scales. The solid blue
curves represent the result of the simultaneous fit described in the text,
with the following components: $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ signal
(dotted black), nonresonant $B^{\pm}\\!\rightarrow K^{+}K^{-}\pi^{\pm}$
background (dashed black), $B^{\pm}\\!\rightarrow\phi K^{\pm}$ signal (dotted
red), nonresonant $B^{\pm}\\!\rightarrow K^{+}K^{-}K^{\pm}$ background (dashed
red), partially-reconstructed $b$-hadron background with (dotted blue) or
without (dashed blue) a true $\phi$ meson, and combinatorial background with
(dotted green) or without (dashed green) a true $\phi$ meson. Normalized
residuals are displayed below the histograms.
The search for $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ decays is performed using
a simultaneous fit to the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ (${\rm
DLL}_{K\pi}<-1$) and $B^{\pm}\\!\rightarrow\phi K^{\pm}$ (${\rm
DLL}_{K\pi}\geq-1$) candidates, dividing the
$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ candidates in four subsamples according
to their $\delta m$ values, each with its set of eight yields. The fit has a
total of $52$ free parameters: $15$ mass shape parameters, $36$ yields, and
the ratio of the total $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ yield to the total
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ yield.
Figure 3 shows the projections of the fitted function superimposed on the
observed mass distributions of the $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$
candidates. The total $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ signal yield is
found to be $19\pm 19$, while the total $B^{\pm}\\!\rightarrow\phi K^{\pm}$
yield is $(3486\pm 76)+(280\pm 25)$ summing the samples of
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ and $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$
candidates. The fitted yield ratio is
$\frac{N(\mbox{$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$})}{N(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})}=(\rm 5.1\,^{+5.3}_{-5.0}(stat)\pm 2.1(syst))\times 10^{-3}\,,$
(12)
where the systematic uncertainty is the quadratic sum of contributions due to
the modelling of the mass shapes ($\pm 2.1\times 10^{-3}$), the fit procedure
($\pm 0.2\times 10^{-3}$), and interference effects between the $\phi$
resonance and a $K^{+}K^{-}$ pair in an S-wave state ($\pm 0.4\times
10^{-3}$). The first contribution is obtained by repeating the fit with the
parameter values of the Crystal Ball and ARGUS functions changed within their
uncertainties, or with an exponential (rather than linear) combinatorial
background model. The dominant effect is due to the 8% uncertainty on the
ratio $\rho$ of the $B^{\pm}\\!\rightarrow\phi K^{\pm}$ mass resolutions in
the two ${\rm DLL}_{K\pi}$ regions. Simulation studies show that the fit
procedure is unbiased, and the statistical precision of this check is assigned
as a systematic uncertainty.
The measurement of the branching fraction ratio is obtained as the ratio
between Eq. 12 and Eq. 3:
$\frac{{\cal B}(\mbox{$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$})}{{\cal
B}(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})}=\rm(6.6\,^{+6.9}_{-6.6}(stat)\pm 2.8(syst))\times 10^{-3}\,.$ (13)
Since the result is not significantly different from zero, we also quote upper
limits from the integral of the likelihood function of this ratio, considering
only the physical (non-negative) region. Including systematic uncertainties we
obtain ${\cal B}(\mbox{$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$})/{\cal
B}(\mbox{$B^{\pm}\\!\rightarrow\phi K^{\pm}$})<0.018\,(0.020)$ at 90% (95%)
CL. Using the current world average ${\cal B}(\mbox{$B^{\pm}\\!\rightarrow\phi
K^{\pm}$})=(8.8\,^{+0.7}_{-0.6})\times 10^{-6}$ [1], we finally obtain
$\displaystyle{\cal B}(\mbox{$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$})$
$\displaystyle=$ $\displaystyle(5.8\,^{+6.1}_{-5.8}\pm 2.5)\times 10^{-8}$
(14) $\displaystyle<$ $\displaystyle 1.5\,(1.8)\times 10^{-7}~{}~{}~{}\mbox{at
90\% (95\%) CL}\,.$ (15)
## 7 Conclusions
The difference in charge asymmetries between the $B^{\pm}\\!\rightarrow\phi
K^{\pm}$ and $B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}K^{\pm}$ decay modes is measured in a sample of $pp$ collisions at
7$\mathrm{\,Te\kern-1.00006ptV}$ centre-of-mass energy, corresponding to an
integrated luminosity of 1.0$\mbox{\,fb}^{-1}$ collected with the LHCb
detector. Using the known value of the
$B^{\pm}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}K^{\pm}$
asymmetry, the $C\\!P$-violating charge asymmetry of
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ decays is determined to be
$\mathcal{A}_{C\\!P}(\mbox{$B^{\pm}\\!\rightarrow\phi K^{\pm}$})=\rm 0.022\pm
0.021(stat)\pm 0.009(syst)$. This result is almost a factor two more precise
than the current world average [1]. It is consistent with both the absence of
$C\\!P$ violation and the Standard Model prediction.
A search for $B^{\pm}\\!\rightarrow\phi\pi^{\pm}$ decays is also performed. No
significant signal is found. Using the known branching fraction of the
$B^{\pm}\\!\rightarrow\phi K^{\pm}$ normalization channel, an upper limit of
${\cal B}(\mbox{$B^{\pm}\\!\rightarrow\phi\pi^{\pm}$})<1.5\,(1.8)\times
10^{-7}$ is set at 90% (95%) confidence level. This improves on the previous
best upper limit [19], while reaching the upper end of the Standard Model
predictions.
## Acknowledgements
We express our gratitude to our colleagues in the CERN accelerator departments
for the excellent performance of the LHC. We thank the technical and
administrative staff at the LHCb institutes. We acknowledge support from CERN
and from the national agencies: CAPES, CNPq, FAPERJ and FINEP (Brazil); NSFC
(China); CNRS/IN2P3 and Region Auvergne (France); BMBF, DFG, HGF and MPG
(Germany); SFI (Ireland); INFN (Italy); FOM and NWO (The Netherlands); SCSR
(Poland); MEN/IFA (Romania); MinES, Rosatom, RFBR and NRC “Kurchatov
Institute” (Russia); MinECo, XuntaGal and GENCAT (Spain); SNSF and SER
(Switzerland); NAS Ukraine (Ukraine); STFC (United Kingdom); NSF (USA). We
also acknowledge the support received from the ERC under FP7. The Tier1
computing centres are supported by IN2P3 (France), KIT and BMBF (Germany),
INFN (Italy), NWO and SURF (The Netherlands), PIC (Spain), GridPP (United
Kingdom). We are thankful for the computing resources put at our disposal by
Yandex LLC (Russia), as well as to the communities behind the multiple open
source software packages that we depend on.
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|
arxiv-papers
| 2013-09-15T09:55:33 |
2024-09-04T02:49:50.909284
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "LHCb collaboration: R. Aaij, B. Adeva, M. Adinolfi, C. Adrover, A.\n Affolder, Z. Ajaltouni, J. Albrecht, F. Alessio, M. Alexander, S. Ali, G.\n Alkhazov, P. Alvarez Cartelle, A.A. Alves Jr, S. Amato, S. Amerio, Y. Amhis,\n L. Anderlini, J. Anderson, R. Andreassen, J.E. Andrews, R.B. Appleby, O.\n Aquines Gutierrez, F. Archilli, A. Artamonov, M. Artuso, E. Aslanides, G.\n Auriemma, M. Baalouch, S. Bachmann, J.J. Back, A. Badalov, C. Baesso, V.\n Balagura, W. Baldini, R.J. Barlow, C. Barschel, S. Barsuk, W. Barter, Th.\n Bauer, A. Bay, J. Beddow, F. Bedeschi, I. Bediaga, S. Belogurov, K. Belous,\n I. Belyaev, E. Ben-Haim, G. Bencivenni, S. Benson, J. Benton, A. Berezhnoy,\n R. Bernet, M.-O. Bettler, M. van Beuzekom, A. Bien, S. Bifani, T. Bird, A.\n Bizzeti, P.M. Bj{\\o}rnstad, T. Blake, F. Blanc, J. Blouw, S. Blusk, V. Bocci,\n A. Bondar, N. Bondar, W. Bonivento, S. Borghi, A. Borgia, T.J.V. Bowcock, E.\n Bowen, C. Bozzi, T. Brambach, J. van den Brand, J. Bressieux, D. Brett, M.\n Britsch, T. Britton, N.H. Brook, H. Brown, A. Bursche, G. Busetto, J.\n Buytaert, S. Cadeddu, O. Callot, M. Calvi, M. Calvo Gomez, A. Camboni, P.\n Campana, D. Campora Perez, A. Carbone, G. Carboni, R. Cardinale, A. Cardini,\n H. Carranza-Mejia, L. Carson, K. Carvalho Akiba, G. Casse, L. Castillo\n Garcia, M. Cattaneo, Ch. Cauet, R. Cenci, M. Charles, Ph. Charpentier, P.\n Chen, S.-F. Cheung, N. Chiapolini, M. Chrzaszcz, K. Ciba, X. Cid Vidal, G.\n Ciezarek, P.E.L. Clarke, M. Clemencic, H.V. Cliff, J. Closier, C. Coca, V.\n Coco, J. Cogan, E. Cogneras, P. Collins, A. Comerma-Montells, A. Contu, A.\n Cook, M. Coombes, S. Coquereau, G. Corti, B. Couturier, G.A. Cowan, D.C.\n Craik, M. Cruz Torres, S. Cunliffe, R. Currie, C. D'Ambrosio, P. David,\n P.N.Y. David, A. Davis, I. De Bonis, K. De Bruyn, S. De Capua, M. De Cian,\n J.M. De Miranda, L. De Paula, W. De Silva, P. De Simone, D. Decamp, M.\n Deckenhoff, L. Del Buono, N. D\\'el\\'eage, D. Derkach, O. Deschamps, F.\n Dettori, A. Di Canto, H. Dijkstra, M. Dogaru, S. Donleavy, F. Dordei, A.\n Dosil Su\\'arez, D. Dossett, A. Dovbnya, F. Dupertuis, P. Durante, R.\n Dzhelyadin, A. Dziurda, A. Dzyuba, S. Easo, U. Egede, V. Egorychev, S.\n Eidelman, D. van Eijk, S. Eisenhardt, U. Eitschberger, R. Ekelhof, L. Eklund,\n I. El Rifai, Ch. Elsasser, A. Falabella, C. F\\\"arber, C. Farinelli, S. Farry,\n D. Ferguson, V. Fernandez Albor, F. Ferreira Rodrigues, M. Ferro-Luzzi, S.\n Filippov, M. Fiore, C. Fitzpatrick, M. Fontana, F. Fontanelli, R. Forty, O.\n Francisco, M. Frank, C. Frei, M. Frosini, E. Furfaro, A. Gallas Torreira, D.\n Galli, M. Gandelman, P. Gandini, Y. Gao, J. Garofoli, P. Garosi, J. Garra\n Tico, L. Garrido, C. Gaspar, R. Gauld, E. Gersabeck, M. Gersabeck, T.\n Gershon, Ph. Ghez, V. Gibson, L. Giubega, V.V. Gligorov, C. G\\\"obel, D.\n Golubkov, A. Golutvin, A. Gomes, P. Gorbounov, H. Gordon, M. Grabalosa\n G\\'andara, R. Graciani Diaz, L.A. Granado Cardoso, E. Graug\\'es, G. Graziani,\n A. Grecu, E. Greening, S. Gregson, P. Griffith, O. Gr\\\"unberg, B. Gui, E.\n Gushchin, Yu. Guz, T. Gys, C. Hadjivasiliou, G. Haefeli, C. Haen, S.C.\n Haines, S. Hall, B. Hamilton, T. Hampson, S. Hansmann-Menzemer, N. Harnew,\n S.T. Harnew, J. Harrison, T. Hartmann, J. He, T. Head, V. Heijne, K.\n Hennessy, P. Henrard, J.A. Hernando Morata, E. van Herwijnen, M. He\\ss, A.\n Hicheur, E. Hicks, D. Hill, M. Hoballah, C. Hombach, W. Hulsbergen, P. Hunt,\n T. Huse, N. Hussain, D. Hutchcroft, D. Hynds, V. Iakovenko, M. Idzik, P.\n Ilten, R. Jacobsson, A. Jaeger, E. Jans, P. Jaton, A. Jawahery, F. Jing, M.\n John, D. Johnson, C.R. Jones, C. Joram, B. Jost, M. Kaballo, S. Kandybei, W.\n Kanso, M. Karacson, T.M. Karbach, I.R. Kenyon, T. Ketel, B. Khanji, O.\n Kochebina, I. Komarov, R.F. Koopman, P. Koppenburg, M. Korolev, A.\n Kozlinskiy, L. Kravchuk, K. Kreplin, M. Kreps, G. Krocker, P. Krokovny, F.\n Kruse, M. Kucharczyk, V. Kudryavtsev, K. Kurek, T. Kvaratskheliya, V.N. La\n Thi, D. Lacarrere, G. Lafferty, A. Lai, D. Lambert, R.W. Lambert, E.\n Lanciotti, G. Lanfranchi, C. Langenbruch, T. Latham, C. Lazzeroni, R. Le Gac,\n J. van Leerdam, J.-P. Lees, R. Lef\\`evre, A. Leflat, J. Lefran\\c{c}ois, S.\n Leo, O. Leroy, T. Lesiak, B. Leverington, Y. Li, L. Li Gioi, M. Liles, R.\n Lindner, C. Linn, B. Liu, G. Liu, S. Lohn, I. Longstaff, J.H. Lopes, N.\n Lopez-March, H. Lu, D. Lucchesi, J. Luisier, H. Luo, O. Lupton, F. Machefert,\n I.V. Machikhiliyan, F. Maciuc, O. Maev, S. Malde, G. Manca, G. Mancinelli, J.\n Maratas, U. Marconi, P. Marino, R. M\\\"arki, J. Marks, G. Martellotti, A.\n Martens, A. Mart\\'in S\\'anchez, M. Martinelli, D. Martinez Santos, D. Martins\n Tostes, A. Martynov, A. Massafferri, R. Matev, Z. Mathe, C. Matteuzzi, E.\n Maurice, A. Mazurov, J. McCarthy, A. McNab, R. McNulty, B. McSkelly, B.\n Meadows, F. Meier, M. Meissner, M. Merk, D.A. Milanes, M.-N. Minard, J.\n Molina Rodriguez, S. Monteil, D. Moran, P. Morawski, A. Mord\\`a, M.J.\n Morello, R. Mountain, I. Mous, F. Muheim, K. M\\\"uller, R. Muresan, B. Muryn,\n B. Muster, P. Naik, T. Nakada, R. Nandakumar, I. Nasteva, M. Needham, S.\n Neubert, N. Neufeld, A.D. Nguyen, T.D. Nguyen, C. Nguyen-Mau, M. Nicol, V.\n Niess, R. Niet, N. Nikitin, T. Nikodem, A. Nomerotski, A. Novoselov, A.\n Oblakowska-Mucha, V. Obraztsov, S. Oggero, S. Ogilvy, O. Okhrimenko, R.\n Oldeman, M. Orlandea, J.M. Otalora Goicochea, P. Owen, A. Oyanguren, B.K.\n Pal, A. Palano, M. Palutan, J. Panman, A. Papanestis, M. Pappagallo, C.\n Parkes, C.J. Parkinson, G. Passaleva, G.D. Patel, M. Patel, G.N. Patrick, C.\n Patrignani, C. Pavel-Nicorescu, A. Pazos Alvarez, A. Pearce, A. Pellegrino,\n G. Penso, M. Pepe Altarelli, S. Perazzini, E. Perez Trigo, A. P\\'erez-Calero\n Yzquierdo, P. Perret, M. Perrin-Terrin, L. Pescatore, E. Pesen, G. Pessina,\n K. Petridis, A. Petrolini, A. Phan, E. Picatoste Olloqui, B. Pietrzyk, T.\n Pila\\v{r}, D. Pinci, S. Playfer, M. Plo Casasus, F. Polci, G. Polok, A.\n Poluektov, E. Polycarpo, A. Popov, D. Popov, B. Popovici, C. Potterat, A.\n Powell, J. Prisciandaro, A. Pritchard, C. Prouve, V. Pugatch, A. Puig\n Navarro, G. Punzi, W. Qian, B. Rachwal, J.H. Rademacker, B.\n Rakotomiaramanana, M.S. Rangel, I. Raniuk, N. Rauschmayr, G. Raven, S.\n Redford, S. Reichert, M.M. Reid, A.C. dos Reis, S. Ricciardi, A. Richards, K.\n Rinnert, V. Rives Molina, D.A. Roa Romero, P. Robbe, D.A. Roberts, A.B.\n Rodrigues, E. Rodrigues, P. Rodriguez Perez, S. Roiser, V. Romanovsky, A.\n Romero Vidal, J. Rouvinet, T. Ruf, F. Ruffini, H. Ruiz, P. Ruiz Valls, G.\n Sabatino, J.J. Saborido Silva, N. Sagidova, P. Sail, B. Saitta, V. Salustino\n Guimaraes, B. Sanmartin Sedes, R. Santacesaria, C. Santamarina Rios, E.\n Santovetti, M. Sapunov, A. Sarti, C. Satriano, A. Satta, M. Savrie, D.\n Savrina, M. Schiller, H. Schindler, M. Schlupp, M. Schmelling, B. Schmidt, O.\n Schneider, A. Schopper, M.-H. Schune, R. Schwemmer, B. Sciascia, A. Sciubba,\n M. Seco, A. Semennikov, K. Senderowska, I. Sepp, N. Serra, J. Serrano, P.\n Seyfert, M. Shapkin, I. Shapoval, Y. Shcheglov, T. Shears, L. Shekhtman, O.\n Shevchenko, V. Shevchenko, A. Shires, R. Silva Coutinho, M. Sirendi, N.\n Skidmore, T. Skwarnicki, N.A. Smith, E. Smith, E. Smith, J. Smith, M. Smith,\n M.D. Sokoloff, F.J.P. Soler, F. Soomro, D. Souza, B. Souza De Paula, B.\n Spaan, A. Sparkes, P. Spradlin, F. Stagni, S. Stahl, O. Steinkamp, S.\n Stevenson, S. Stoica, S. Stone, B. Storaci, M. Straticiuc, U. Straumann, V.K.\n Subbiah, L. Sun, W. Sutcliffe, S. Swientek, V. Syropoulos, M. Szczekowski, P.\n Szczypka, D. Szilard, T. Szumlak, S. T'Jampens, M. Teklishyn, E. Teodorescu,\n F. Teubert, C. Thomas, E. Thomas, J. van Tilburg, V. Tisserand, M. Tobin, S.\n Tolk, D. Tonelli, S. Topp-Joergensen, N. Torr, E. Tournefier, S. Tourneur,\n M.T. Tran, M. Tresch, A. Tsaregorodtsev, P. Tsopelas, N. Tuning, M. Ubeda\n Garcia, A. Ukleja, A. Ustyuzhanin, U. Uwer, V. Vagnoni, G. Valenti, A.\n Vallier, R. Vazquez Gomez, P. Vazquez Regueiro, C. V\\'azquez Sierra, S.\n Vecchi, J.J. Velthuis, M. Veltri, G. Veneziano, M. Vesterinen, B. Viaud, D.\n Vieira, X. Vilasis-Cardona, A. Vollhardt, D. Volyanskyy, D. Voong, A.\n Vorobyev, V. Vorobyev, C. Vo\\ss, H. Voss, R. Waldi, C. Wallace, R. Wallace,\n S. Wandernoth, J. Wang, D.R. Ward, N.K. Watson, A.D. Webber, D. Websdale, M.\n Whitehead, J. Wicht, J. Wiechczynski, D. Wiedner, L. Wiggers, G. Wilkinson,\n M.P. Williams, M. Williams, F.F. Wilson, J. Wimberley, J. Wishahi, W.\n Wislicki, M. Witek, G. Wormser, S.A. Wotton, S. Wright, S. Wu, K. Wyllie, Y.\n Xie, Z. Xing, Z. Yang, X. Yuan, O. Yushchenko, M. Zangoli, M. Zavertyaev, F.\n Zhang, L. Zhang, W.C. Zhang, Y. Zhang, A. Zhelezov, A. Zhokhov, L. Zhong, A.\n Zvyagin",
"submitter": "Jessica Prisciandaro",
"url": "https://arxiv.org/abs/1309.3742"
}
|
1309.3744
|
# CaMn2Sb${}_{2}:$ a fully frustrated classical magnetic system
I. I. Mazin Code 6393, Naval Research Laboratory, Washington, DC 20375, USA
Code 6393, Naval Research Laboratory, Washington, DC 20375, USA
(8/9/13)
###### Abstract
We show by means of density functional calculations that CaMn2Sb2 is very
close to a mean-field critical point known for the classical Heisenberg model
on the honeycomb lattice. Three entirely different long range ordered magnetic
phases become degenerate at this point: the Neel phase and two different
spiral phases. We speculate that the unusual physical properties of this
compound, observed in recent experiments, in particular the enigmatic
intermediate temperature phase, are due to this proximity.
Introduction: CaMn2Sb2 is an interesting implementation of a nearly classical
spin system on a honeycomb lattice. Such systems have been attracting the
attention of researchers since the seminal paper of Rastelli $et$ $al$ from
1979Rastelli . The basic model involves the nearest neighbor (nn) interaction
$J_{1}$ and the second neighbor interaction $J_{2}.$ In the classical limit,
the ground state of this $J_{1}-J_{2}$ model, for antiferromagnetic $J_{1,2}$,
is the Neel phase where all nearest neighbor bonds are fully
antiferromagnetic, if $J_{2}/J_{1}<1/6.$ Interestingly, for
$1/6<J_{2}/J_{1}<1/2,$ there are two degenerate solutions, corresponding to
two different spiral phases (in one of them the two spins in one unit cell are
always antiparallel, and the spiral propagation vector is perpendicular to a
nn bond, in the other the two sublattices are rotated by a particular angle
with respect to each other, and the spiral vector is parallel to a bond). At
$J_{2}/J_{1}>1/2$ the ground state is the stripe phase with alternating FM
pairs (Fig. 1). Both critical points, $J_{2}/J_{1}=1/2$ and $J_{2}/J_{1}=1/6$
are triple points: in case there is an additional parameter, for instance the
third neighbor exchange $J_{3},$ three phases meet at these points: the Neel
phase and the two spiral phases at $J_{2}/J_{1}=1/6,$ and the stripe phase and
the two spiral phases at $J_{2}/J_{1}=1/2.$ Whenever a third parameter is
added, be it the third neighbor exchangeRastelli $J_{3}$, uniaxial
anisotropyRastelli or biquadratic couplingI $K,$ complex phase diagrams
arise, with the four phases described above, and additional phases such as
zigzag antiferromagnetism.
Figure 1: Four antiferromagnetic patterns discussed in the text. Note that
the spiral A coincides with the phase N at $\theta=0$ and spiral B with the
phase N N at $\phi=\pi$ and $\theta=0$, and with S at $\phi=0$ and
$\theta=\pi$. At the same time the two spirals cannot be continuously
transformed into each other except in this limiting cases.
Figure 2: (color online) Crystal structure of CaMn2Sb2, showing examples of
the nearest neighbors ($J_{1}$, red), 2nd neighbors ($J_{2}$, green) and 3rd
neighbors ($J_{3}$), blue, inside a MnSb layer, as well as the nearest
($J_{z1}$, yellow) and 2nd ($J_{z2}$, black) interplanar neighbors. Mn atoms
are shown by gray spheres, Sb by cyan ones and Ca by green ones.
The crystal structure of CaMn2Sb2 is shown in Fig. 2. The Mn sublattice
consists of honeycomb layers in which every other atom is shifted
perpendicular to the plane. This interesting geometry leads to a rather short
Mn-Mn distance, 3.18 Å, with substantial direct overlap between the Mn d
orbitals and as a result sizable direct antiferromagnetic exchange. In
addition, there are two superexchange paths available. One connects the
nearest neghbor Mns, with an Mn-Sb-Mn angle of 70${}^{\circ},$ and the other
connects the second neighbors, with an angle of 108${}^{\circ}.$ Mn in this
compound has valency 2+, and therefore in the high spin state it would have a
magnetic moment of 5 $\mu_{B},$ reduced by hybridization and fluctuations.
Experimentally at low temperature the $J_{1}$-driven Neel phase was found,
with a magnetic moment of 2.8-3.4 $\mu_{B}$. At $T=85$ K this phase gives rise
to another phase of unknown origin, which experimentally resembles weak
ferromagnetism, but no long range antiferromagnetic order was detected.
Finally, at $T\gtrsim 200$ K the material becomes paramagnetic, exhibiting at
$T\gtrsim 300$ K a Curie-Weiss behavior with an effective moment $M_{CW}=1.4$
$\mu_{B}$. The unusual character of the intermediate temperature phase, as
well as the very low Curie-Weiss moment, suggest that frustration
characteristic of honeycomb magnetic models may play a role. The low-
temperature phase shows an activation transport behavior with an activation
gap $\sim 40$ meV, while the intermediate temperature phase exhibits a strong
increase (up to a factor of 100) of the resistivity.Aronson
We have performed first principle calculations of the electronic and magnetic
properties of CaMn2Sb2. We found Mn to be in the high spin state and
$d({5\uparrow})$ configuration. We also found a small gap consistent with the
experiment, and magnetic interactions dominated by the nearest neighbor
exchange. Most interestingly, we found that the second neighbor exchange is
4-6 times smaller than the nearest neighbor one, while the third neighbor
exchange and the biquadratic exchange are very small, and the magnetic
anisotropy is of the easy-plane type. In this regime classical spins on the
honeycomb lattice are highly frustrated, with two spiral and two collinear
phases nearly degenerate. We speculate that this frustration is the cause of
the unusual magnetic phase diagram.
Calculations. The calculations in this paper were performed with the Linear
Augmented Plane Wave code WIEN2kWIEN2K , using the following crystallography:
symmetry group #164, $P\bar{3}m1,$ $a=$4.522 Å, $c=$7.458 Å, $z_{Ca}=0,$
$z_{Mn}=0.3784,$ $z_{Sb}=0.7487,$ and a Generalized Gradient Approximation for
the exchange-correlation potential. In agreement with previous
calculationsAronson we found that Mn in strongly polarized, and Mn moments
are well localized, as demonstrated by the fact that all magnetic
configurations converge to about the same magnetic moment, $\sim 4$ $\mu_{B},$
and the calculated exchange energies are much smaller than the magnetization
energy. The lowest energy among various collinear states has the Neel state,
in which the nn Mn have opposite spins.
Figure 3: (color online) Calculated band structure of CaMn2Sb2 in the
antiferromagnetic Neel phase. The symbol size is proportional to the Mn
character (spin-up in blue, below the Fermi level, and spin-down in red, above
the Fermi level. Figure 4: The calculated dielectric function of CaMn2Sb2.
The inset shows the corresponding transmission coefficient (the thickness of
the slab was chosen to roughly match the amplitude of the transmission
coefficient in Ref. Aronson i, and a uniform week absorption was added to
account for the experimentally observed in-gap states).
The band structure and density of states in this configuration are shown in
Fig. 3. One can clearly see that all five $d$ orbitals of Mn are fully
polarized, while the reduction of the moment from 5 to 4 $\mu_{B}$ is due to
hybridization. There is also an indirect excitation gap (of the spin-flip
nature), about 50 meV, consistent with the experimentally measured one. The
calculated optical gap (minimal direct transition energy) is about 0.7 meV,
but the absorption at this energy starts relatively slow (Fig. 4), owing to
the fact that the top of the valence band and the bottom of the conductivity
band are formed in the different spin channels (in a ferromagnetic case, these
transtions would have been optically forbidden; in an antiferromagnetic one,
they are allowed but weak). The main absorption edge is located at 2 eV.
Compared to experimentAronson , the calculated transmission drops at a lower
frequency. However, the optical experiments reported in Ref. Aronson should
be considered as preliminary. One has to wait for more state-of-the-art
optical data, preferably for the conductivity, for a full comparison.
In order to understand better magnetic interactions in CaMn2Sb${}_{2},$ we
have calculated the total energies of 10 different magnetic configurations,
and mapped the results onto the model Heisenberg that includes 1st, 2nd and
3rd nearest neighbors inside the MnSb layers ($J_{1},$ $J_{2},$ and $J_{3})$,
and 1st and 2nd neighbors between the layers ($J_{z1}$ and $J_{z2}).$ The
corresponding distances and values are listed in the Table. Note that the
definition of $J$ in this paper is such that the pair interaction is
$J\sigma_{i}\sigma_{j},$ where $\sigma=\pm 1.$
Table 1: Calculated exchange parameters and the corresponding Mn-Mn bond lengths. See Fig.2 for graphical examples. | $J_{1}$ | $J_{2}$ | $J_{3}$ | $J_{z1}$ | $J_{z2}$
---|---|---|---|---|---
$d$ (Å) | 3.179 | 4.522 | 5.528 | 6.219 | 7.458
$J$ (meV) | 91 | 22 | 5 | 11 | 3
Since the number of energy differences (9) is larger than the number of fitted
exchange constants, one can estimate the accuracy of the fitting from the
standard fit error. This is on the order of 7–8 meV. Note that on the mean
field level $J_{z}$s simply renormalize the planar interactions, in other
words, assuming perfect interplanar ordering, one can map the 3D problem onto
2D with $J_{1eff}=J_{1}+J_{z1}=102$ meV.
We have also looked for other possible deviations from the Heisenberg model:
the biquadratic interaction, which plays an extraordinary important role in
Fe-based superconductors, and single-site anisotropy. The former appears to be
zero within computational accuracy, in other words, for all possible angles
between the Mn spins in the same unit cell the energy is perfectly
cosinusoidal, $E=const+0.5J_{1}\cos\alpha.$ The latter (magnetic anisotropy)
is zero in plane (although formally hexagonal symmetry allows for magnetic
anisotropy when the field is rotated by 30${}^{o},$ for all practical purposes
this effect is negligible), and about 2–3 meV/Fe otherwise, with the direction
perpendicular to the plane being the hard axis.
Discussion. The fact that essentially any magnetic pattern can be stabilized
in the calculations, and that the energy cost of suppressing magnetism
entirely (0.8 eV/Fe) is much larger than the exchange constants indicates that
CaMn2Sb2 should be considered a local moment system, and overall one expects
Mn to be subject to considerable Hubbard correlations. The fact that the
nearest-neighbor exchange follows the Heisenberg formula nearly exactly also
speaks in favor of classical superexchange and the Hubbard model. Yet DFT
calculations reproduce the excitation gap well, in fact, better than dynamic
mean field calculationsHaule . The reasons for such an unexpected success of
the DFT are unclear st the moment. Given the sharp increase of the resistivity
in the other, intermediate temperature phase, one may think that the Neel
phase is affected by some cancellation of errors, not operative in the other
phases, which leads to an effective increase of the excitation gap in the
intermediate phase.
In the classical phase diagram of an antiferromagnetic $J_{1},J_{2}$ 2D
honeycomb layer there are five planar phases: The Neel (N) phase, where all
nearest neighbors order antiferromagnetically, the “stripy” (S) phase where
all bond along one direction are ferromagnetic, and forming double stripes
order antiferromagnetically, and two spiral phases (A and B) described in the
introduction. For $J_{2}/J_{1}<1/6$ the ground state is N, for
$1/6<J_{2}/J_{1}<1$ the ground state is degenerate between phases A and B, and
for the larger $J_{2}$ the ground state is S. Adding $J_{3}$ creates a
complicated phase diagram, with a new “zigzag” phase emerging at large
$J_{3},$ which is shown in Fig. 5, where we also for completeness show the
phase diagram $J_{1},J_{2},K$ with a biquadratic interaction, not published
before. In this phase diagram, the calculated interactions correspond to the
point $J_{2}/J_{1eff}\sim 0.21$, $J_{3}/J_{1eff}\sim 0.05,$ but the error in
these numbers is close to $\pm 0.08.$ In other words, according to the
calculations the material is extremely close to a highly frustrated critical
point where three different phases are degenerate, the Neel phase and two
qualitatively different spiral phases, $J_{2}/J_{1}=1/6,$ $J_{3}=0.$
Figure 5: Phase diagram of the classical $J_{1}--J_{2}--J_{3}$ Heisenberg
model. The phase boundaries are given by the lines
$y=(x-\sqrt{x^{2}+2(x-1/2)^{2}})/2$, $y=(1/2+3x-\sqrt{51/12-5x+9x^{2}})/4$,
$y=(6x-1)/4$, $y=(2x+1)/4$, $y=0$, $x=1/2$Rastelli . The square shows the
region in the phase diagram where, within computational accuracy, the actual
system is located. The insert shows the phase diagram for the $J_{1}-J_{2}-K$
model, where K is the biquadratic coupling. The phase boundaries are:
$y=(3x-\sqrt{1-2x+9x^{2}})/2$, $y=\sqrt{x-1/4}-x$, $y=(6x-1)/2$, $y=0$I .
The thermodynamics of the classical Heisenberg model on the honeycomb lattice
has not, to our knowledge, been systematically studied and this subject is
beyond the scope of the current paper. It is however likely that the phase
boundaries between the zero-temperature phase shift as the entropy is
included, so it is tempting to ascribe the transformation at 85 K to a phase
transition between the Neel and one of the two spiral phases (more likely the
spiral A). Very weak net ferromagnetism is then due to relativistic effects,
which can be estimated by the ratio between the magnetic anisotropy energy
(2-3 meV) and magnetization energy (800 meV), which is $\sim 3\times 10^{-3},$
to be compared with the experimental ferromagnetic moment of $7\times 10^{-3}$
$\mu_{B}\approx 3\times 10^{-3}$ $M_{Mn}.$
No long range order has been observed between 85 and 210 K in neutron
scattering; this suggests that the spirals break after relatively short
distance and the emerging spiral chunks are randomly distributed over the
three equivalent crystallographic directions. Such a state would be very
similar to the “magnetic liquid” state in MnSiReznik , also called a
“cholesteric” phase, or a “blue” phase in the original paper. This state is
observed in MnSi near the phase transition into a long-range ordered spiral
phase, on the high-temperature (high-pressure) side of the phase transition.
Neutron scattering in this phase reveals no long range order, but well-defined
spirals of considerable length (several hundred Å), propagating in all
crystallographically equivalent directions with equal weight. It is still not
clear whether the spirals in the “blue phase” of MnSi form domains (possibly
dynamic) or are meandering around, periodically switching directions. Both
options are open as well in the intermediate temperature phase of
CaMn2Sb${}_{2}.$
Another corollary of proximity to the critical point at $J_{2}/J_{1}=1/6$ is
that unusual low energy magnetic excitations should be present in the low-
temperature Neel phase. More detailed experimental spectroscopic studies are
highly desirable.
To summarize, we have shown that CaMn2Sb2 is extremely close (within
computational accuracy) to a critical point where three entirely different
magnetic phases are degenerate on the mean field level. We suggest that
unusual properties of this compound, obtained in recent experiments, are
related to this unique proximity.
The author acknowledges discussions with Girsh Blumberg and Meigan Aronson,
and funding from the Office of Naval Research (ONR) through the Naval Research
Laboratory’s Basic Research Program, and from the Alexander von Humboldt
Foundation.
## References
* (1) E. Rastelli, A. Tassi, and L. Reatto, Physica B 97, 1 (1979).
* (2) J.B. Fouet, P. Sindzingre, and C. Lhuillier, Eur. Phys. J. B 20, 241 (2001).
* (3) I.I. Mazin, unpublished.
* (4) J. W. Simonson, G. J. Smith, K. Post, M. Pezzoli, J. J. Kistner-Morris, D. E. McNally, J. E. Hassinger, C. S. Nelson, G. Kotliar, D. N. Basov, and M. C. Aronson, Phys. Rev. B 86, 184430 (2012).
* (5) P. Blaha, K. Schwarz, G. K. H. Madsen, D. Kvasnicka and J. Luitz, WIEN2k, An Augmented Plane Wave + Local Orbitals Program for Calculating Crystal Properties (K. Schwarz, Techn. Universit$\ddot{a}$t Wien, Austria, 2001), ISBN 3-9501031-1-2.
* (6) K. Haule, unpublished
* (7) C. Pfleiderer, D. Reznik, L. Pintschovius, H. v. Löhneysen1, M. Garst and A. Rosch, Nature 427, 227 (2004).
|
arxiv-papers
| 2013-09-15T10:23:00 |
2024-09-04T02:49:50.920300
|
{
"license": "Public Domain",
"authors": "I. I. Mazin",
"submitter": "Igor Mazin",
"url": "https://arxiv.org/abs/1309.3744"
}
|
1309.3752
|
# Novel Repair-by-Transfer Codes and Systematic Exact-MBR Codes with Lower
Complexities and Smaller Field Sizes
Sian-Jheng Lin and Wei-Ho Chung Authors are with the Research Center for
Information Technology Innovation, Academia Sinica, Taipei City, Taiwan.
(e-mail: [email protected]; [email protected])
###### Abstract
The $(n,k,d)$ regenerating code is a class of $(n,k)$ erasure codes with the
capability to recover a lost code fragment from other $d$ existing code
fragments. This paper concentrates on the design of exact regenerating codes
at Minimum Bandwidth Regenerating (MBR) points. For $d=n-1$, a class of
$(n,k,d=n-1)$ Exact-MBR codes, termed as repair-by-transfer codes, have been
developed in prior work to avoid arithmetic operations in node repairing
process. The first result of this paper presents a new class of repair-by-
transfer codes via congruent transformations. As compared with the prior work,
the advantages of the proposed codes include: i). The minimum of the finite
field size is significantly reduced from $n\choose 2$ to $n$. ii). The
encoding complexity is decreased from $n^{4}$ to $n^{3}$. As shown in
simulations, the proposed repair-by-transfer codes have lower computational
overhead when $n$ is greater than a specific constant. The second result of
this paper presents a new form of coding matrix for product-matrix Exact-MBR
codes. The proposed coding matrix includes a number of advantages: i). The
minimum of the finite field size is reduced from $n-k+d$ to $n$. ii). The fast
Reed-Solomon erasure coding algorithms can be applied on the Exact-MBR codes
to reduce the time complexities.
###### Index Terms:
Distributed storage, maximum-distance-separable (MDS) codes, partial
downloading, Reed-Solomon codes, repair-by-transfer.
## 1 Introduction
In a distributed storage system, the source data (message) is dispersed across
nodes in the network, and a data collector (DC) can retrieve the whole source
data by accessing a subset of the nodes. To tradeoff between the storage
efficiency and the system reliability, the erasure codes, such as maximum-
distance-separable (MDS) codes [1], random linear codes [2] or fountain codes
[3, 4], are usually adopted as the base of data format in distributed storage
systems [5, 6, 27]. For an unstable network, the nodes may frequently join and
depart. When a node departs or crashes, the system manager will place a
replacement node in the distributed storage network to replace the
functionality of the failed node. Suppose the replacement node does not store
any information about the data (code fragments) stored in the prior failed
node. To reconstruct the data, the replacement node broadcasts a request to a
subset of other helper nodes, and those helper nodes reply with the requisite
information to the replacement node. If the distributed storage systems is
based on conventional Reed-Solomon (RS) codes, an intuitive method is to
reconstruct the entire source data in the replacement node, and then extract
the desired code fragment from the source data. By such method, the total
amount of downloaded symbols is not less than the size of whole source data.
However, as the size of data stored in a single node is much smaller than the
entire source data, it is possible to design a new class of storage codes to
reduce the amount of downloaded symbols in node-repairing process. The new
class of storage codes, termed as regenerating codes, is introduced by the
pioneer paper [7].
### 1.1 Coding system description
In this paper, the regenerating code over $GF(q)$ is associated with a set of
parameters $\\{n,k,d,\alpha,\beta,B\\}$ elaborated in the following. The value
$B$ is the number of source symbols over $GF(q)$ to be encoded. The $n$ is the
number of produced code fragments, which will be respectively stored in $n$
network nodes. The $\alpha$ is the number of symbols of a code fragment. In
data reconstruction process, the DC individually downloads $\alpha$ symbols
from each of a subset of $k$ nodes to reconstruct the message. In the node-
repairing process, the replacement node individually downloads $\beta$ symbols
from each of a subset of $d$ integrity nodes to rebuild the code fragment.
Those parameters $\\{n,k,d\\}$ follows the inequality
$k\leq d\leq n-1.$
The theoretical bound of storage-bandwidth trade-off have been given by [8]
based on the cut-set bound of network coding:
$B\leq\sum_{i=0}^{k-1}min\\{\alpha,(d-i)\beta\\}.$ (1)
By the theoretical bound (1), two extreme points on the storage-bandwidth
trade-off have been adequately investigated in prior works. The first extreme
point, termed as minimum storage regeneration (MSR) point, is firstly to
minimize the $\alpha$ and then minimize the $\beta$. The parameter
configuration is
$\begin{array}[]{l}\alpha=B/k;\\\ \beta=B/\left(k(d-k+1)\right).\end{array}$
(2)
The second extreme point, termed as the minimum bandwidth regenerating (MBR)
point, is firstly to minimize the $\beta$, and then minimize the $\alpha$. The
parameter configuration is
$\begin{array}[]{l}\beta=2B/\left(k(2d-k+1)\right);\\\
\alpha=d\beta.\end{array}$ (3)
By the so-called data striping technique [9], the regenerating codes at
$\beta=1$ can be used to construct the regenerating codes for any $\beta$.
Thus, here in after, we focus on the design of regenerating codes at the
$\beta=1$ MBR points, and the corresponding parameter configuration is
$\alpha=d,\beta=1\textrm{, and }B={k+1\choose 2}+k(d-k).$ (4)
In the node-regenerating process, if the restored fragment is always the same
with the fragment in the prior failed node, this property is called the exact
regeneration. This is in contrast to the functional regeneration without
imposing restrictions on the content of the stored fragment. Practically, the
exact regeneration is a good property to simplify the hardware and software
designs for distributed storage systems. However, the non-existence of exact
regeneration codes at the interior points on the storage-bandwidth trade-off
curve have been proved [10]. In this paper, the abbreviations ”Exact-MSR” and
”Exact-MBR” respectively indicate the regenerating codes at MSR and MBR points
with the exact regeneration property.
### 1.2 Definitions of terminologies
#### 1.2.1 Systematic regenerating codes
The [9] defines the systematic regenerating code as a class of regenerating
code whose $B$ message symbols appear on a certain set of $k$ systematic code
fragments. The nodes storing those systematic fragments are termed as the
systematic nodes. A major work of this paper is to construct the systematic
regenerating codes at MBR points. Systematic codes are useful in data
reconstruction: If the DC can download those systematic code fragments, the DC
can directly obtain the corresponding pieces of source data without any
computational cost. This is a good property for practical systems.
#### 1.2.2 Repair-by-transfer codes
In the node-repairing process, the replacement node broadcasts a request to a
subset of helper nodes, and each helper node returns certain number of
responding symbols to the replacement node. In general, each helper node
should compute the responding symbols via a function of the fragment stored in
the node. The repair-by-transfer codes are a class of distributed storage
codes where each helper node simply needs to pass a portion of the stored
fragment without any arithmetic operations. The repair-by-transfer codes are
particularly beneficial to the unstable network environment with frequent
occurrence of the node regenerations. A repair-by-transfer code at
$(n,k,d=n-1)$ Exact-MBR case is proposed by Shah et al. [10], and the non-
existence of other cases $d<n-1$ is shown in [19]. The details [10] are
introduced in Section 3. Furthermore, the generalized form of [10] is
presented in [20, 21]. A system implementation for $k=n-1$ and $k=n-2$ is
demonstrated by Hu et al. [26]. A objective of this paper is to construct the
$(n,k,d=n-1)$ Repair-by-transfer codes with smaller finite fields and lower
computational costs. By assigning $d=n-1$ to the (4), the parameters for
$(n,k,d=n-1)$ repair-by-transfer codes are
$\alpha=d=n-1,\beta=1\textrm{, and }B=(n-1)k-{k\choose 2}.$ (5)
#### 1.2.3 Partial downloading scheme
By the MBR data reconstruction process in [9], the DC should download the
whole data stored in the set of connected nodes. To reduce the total amount of
downloaded symbols, Gong and Wang [18] present a data decoding algorithm,
termed as partial downloading scheme, on the non-systematic Exact-MBR codes
[9]. By the partial downloading scheme, the DC can download a partial portion
of code fragment from each connected node. The partial downloading scheme is
useful to mitigate the network congestion. Thus, the partial downloading
schemes are also developed on the proposed repair-by-transfer codes and Exact-
MBR codes.
### 1.3 Previous works
The exact regenerating codes at MSR and MBR points have been proposed in
recent years. For Exact-MSR codes, the [13] discovers the code constructions
at $(n=4,k=2,d=3)$ and ($n=5,k=3,d=4)$ via computer searching. The [14]
presents the Exact-MSR codes for $d=n-1\geq 2k-1$ based on interference
alignment technique. The non-existence of Exact-MSR code for $d<2k-3$ with
$\beta=1$ is shown in [14]. The [15, 16] have shown the existence of exact-MSR
codes for all $(n,k,d)$, while the size of message approaches infinity. By
interference alignment technique, the [17] describes the Exact-MSR codes for
the following cases: i) $k/n\leq 1/2,d\geq 2k-1$; and ii) $k\leq 3$. Rashmi et
al. [9] present an construction for $(n,k,d\geq 2k-2)$ Exact-MSR codes via a
product matrix framework. In Exact-MBR codes, the [10] presents the
$(n=d+1,k,d)$ Exact-MBR codes with no arithmetic operations in node
regeneration process, and the [9] presents the constructions for all feasible
$(n,k,d)$ Exact-MBR codes. Furthermore, the cooperative repair codes [11, 12]
are the generalized version of regenerating codes to address multiple node
failures.
### 1.4 Results and organizations of the paper
In this paper, we developed two classes of Exact-MBR codes. The first result
is the repair-by-transfer code at $(n,k,d=n-1)$ Exact-MBR points via the
congruences of skew-symmetric matrices. The systematic version and the partial
downloading scheme are also proposed. The second result is the systematic
version of Exact-MBR code for all feasible values of $(n,k,d)$ based on the
framework defined by [9]. We design a new encoding matrix for systematic
Exact-MBR code, and the partial downloading scheme are also proposed. To
emphasize the contributions of the paper, Section 5 shows the comparisons of
the proposed codes with the previous works.
Notations and conventions are declared as follows. Throughout this paper, the
operations and symbols are drawn from the field $GF(q)$. For a vector $x$, the
underlined notation as $\underline{x}$ represents a row vector, and the over-
lined notation as $\overline{x}$ represents a column vector. The $x[i]$
denotes the $i$-th element of the vector $x$. For a matrix $X$, the $X[i,j]$
denotes the entry at $i$-th row and $j$-th column. For a matrix (vector) $X$,
the superscript ${}^{\prime}t^{\prime}$ on a matrix (vector) $X^{t}$ denotes
the transpose of this $X$. The $I_{k}$ represents a $k\times k$ identity
matrix.
The rest of this paper is organized as follows. Section 2 reviews the previous
works, such as repair-by-transfer codes and Exact-MBR codes. Section 3
presents the new class of repair-by-transfer codes. Section 4 presents the
proposed systematic Exact-MBR codes based on partially systematic Reed-Solomon
(PSRS) codes. Another construction approach is placed in Appendix. The
comparisons and discussions are placed in Section 5. Section 6 concludes this
paper.
## 2 Previous works
This section reviews a number of related works, such as repair-by transfer
codes [10], Exact-MBR codes [9], and partial downloading scheme [18].
### 2.1 Repair-by-transfer codes [10]
Figure 1: Graphical representation of the $(n=5,k=3,d=4)$ repair-by-transfer
code proposed by [10].
This subsection briefly introduces the $(n,k,d=n-1)$ repair-by-transfer codes
[10] by a simple example $(n=5,k=3,d=4)$ shown in Figure 1. In beginning, the
$B=9$ message symbols are encoded with $(N,K)=\left({n\choose 2},B\right)$ MDS
codes, to generate ${n\choose 2}$ code packets. Then each code packet is
stored in two distinct nodes. The assignment rule can be visualized with a
complete graph of $n$ vertices. As shown in Figure 1, each vertex is
recognized as an individual node, and each edge corresponds to a distinct code
packet. Each node (vertex) stores the $n-1$ code packets linked to this node.
The node regeneration is very simple. If one node fails, the lost $n-1$ code
packets in this node can be directly downloaded from each of other $n-1$
nodes. To reconstruct the data, a DC download the code packets from $k$ nodes.
It can be shown that the DC accesses a total of $B$ distinct code packets, so
the message symbols can be reconstructed via the $\left({n\choose 2},B\right)$
MDS decoding. The [10] suggests that the doubly extended RS codes can be
chosen as the $(N,K)$ MDS coding technique, and the minimal field size is
${n\choose 2}\leq N+1$.
### 2.2 Exact-MBR codes [9] and partial downloading scheme [18]
This section reviews the Exact-MBR codes [9] at (4) through product-matrix
framework. In code constructions, the $B$ message symbols are formed as a
$d\times d$ message matrix $M$, which is then multiplied by an $n\times d$
encoding matrix $\Psi$, resulting in an $n\times d$
$C=\Psi M.$ (6)
code matrix. Let $\underline{c}_{i}^{t}$ denote the $i$-th row of $C$, for
$1\leq i\leq n$. The $\underline{c}_{i}^{t}$ is computed through
$\underline{c}_{i}^{t}=\underline{\psi}_{i}^{t}M,$ (7)
where the $\underline{\psi}_{i}^{t}$ denotes the $i$-th row of $\Psi$. Each
$\underline{c}_{i}^{t}$ is then stored in a network node with index $i$.
The message matrix $M$ is expressed as
$M=\begin{bmatrix}S&T\\\ T^{t}&\mathbf{0}\end{bmatrix},$ (8)
where the $\mathbf{0}$ denotes a $(d-k)\times(d-k)$ zero matrix, the $T$ is a
$k\times(d-k)$ matrix filled with $k(d-k)$ distinct message symbols, and the
$S$ is a $k\times k$ symmetric matrix determined by ${k+1\choose 2}$ message
symbols. The upper triangular part of $S$ is filled with the message symbols,
and other entries assign the corresponding values such that the symmetry
holds. Then, the encoding matrix
$\Psi=\begin{bmatrix}\Phi&\Delta\end{bmatrix}$ (9)
is the concatenation of a $n\times k$ matrix $\Phi$ with a $n\times(d-k)$
matrix $\Delta$. The coding matrix is chosen in such a way that:
i) Any $d$ rows of $\Psi$ are linearly independent;
ii) Any $k$ rows of $\Phi$ are linearly independent.
For the non-systematic case, a feasible form of $\Psi$ is a Vandermonde matrix
[9].
#### 2.2.1 Node-repairing process
Suppose the node $f$ fails, and a replacement node is placed in the network to
replace the functionality of the failure node. To reconstruct the code
fragment (7) in the failure node, the replacement node connects to a subset of
$d$ helper nodes $\\{h_{1},h_{2},\ldots,h_{d}\\}$. Then each helper node
$h_{j}$ computes the scalar value
$\upsilon_{h_{j}}=\underline{c}_{h_{j}}^{t}\underline{\psi}_{f},$ (10)
and passes this value on to the replacement node. Thus, the replacement node
gather $d$ downloaded symbols expressed as a $d$-element column vector
$\Upsilon_{\mathrm{repair}}=[\upsilon_{h_{1}},\upsilon_{h_{2}},\ldots,\upsilon_{h_{d}}]^{t}$.
By definition, the $\Upsilon_{\mathrm{repair}}$ possesses the equality
$\Upsilon_{\mathrm{repair}}=C_{\mathrm{repair}}\underline{\psi}_{f}=\Psi_{\mathrm{repair}}M\underline{\psi}_{f},$
(11)
where the $C_{\mathrm{repair}}$ is a $d\times\alpha$ matrix consisting of $d$
rows
$\\{\underline{c}_{h_{1}}^{t},\underline{c}_{h_{2}}^{t},\ldots,\underline{c}_{h_{d}}^{t}\\}$
taken from the $C$, and the $\Psi_{\mathrm{repair}}$ is a $d\times d$ matrix
consisting of $d$ corresponding encoding rows
$\\{\underline{\psi}_{h_{1}}^{t}\underline{\psi}_{h_{2}}^{t},\ldots,\underline{\psi}_{h_{d}}^{t}\\}$.
As the $\Psi_{\mathrm{repair}}$ is invertible by the first condition of the
MBR encoding matrix, the decoding formula is formulated as
$\Psi_{\mathrm{repair}}^{-1}\times\Upsilon_{\mathrm{repair}}=M\underline{\psi}_{f}=\underline{c}_{f},$
(12)
which is the transpose of the desired fragment $\underline{c}_{f}^{t}$.
#### 2.2.2 Data reconstruction process
To reconstruct the message, the DC connects to $k$ active nodes
$\\{i_{1},i_{2},\ldots,i_{k}\\}$ and then downloads
$\\{\underline{c}_{i_{1}}^{t},\underline{c}_{i_{2}}^{t},\ldots,\underline{c}_{i_{k}}^{t}\\}$
from those connected nodes. The $k$ rows
$\\{\underline{c}_{i_{1}}^{t},\ldots,\underline{c}_{i_{k}}^{t}\\}$ are
formulated as a $k\times\alpha$ matrix $C_{\mathrm{DC}}$ following the order
$[g_{1},\ldots,g_{k}]$. That is, each $\underline{c}_{i_{j}}^{t}$ is placed at
the $g_{j}$-th row of the matrix $C_{\mathrm{DC}}$. In many cases, the
sequence $[g_{1},\ldots,g_{k}]$ can be defined as a monotonically increasing
sequence $g_{i}=i$, $1\leq i\leq k$. However, the proposed partial decoding
scheme, addressed in Sec. 4.3, requires that systematic codeword fragments
should be placed at a specific row of $C_{\mathrm{DC}}$. Based on above
definitions, the DC accesses $k$ vectors expressed as
$C_{\mathrm{DC}}=\Psi_{\mathrm{DC}}M,$ (13)
where the $\Psi_{\mathrm{DC}}$ denotes a $k\times d$ matrix consisting of $k$
corresponding encoding rows
$\\{\underline{\psi}_{i_{1}}^{t},\underline{\psi}_{i_{2}}^{t},\ldots,\underline{\psi}_{i_{k}}^{t}\\}$.
By definition (9), the $k\times d$ matrix $\Psi_{\mathrm{DC}}$ can be
represented as the concatenation of two sub-matrices, given by
$\Psi_{\mathrm{DC}}=\begin{bmatrix}\Phi_{\mathrm{DC}}&\Delta_{\mathrm{DC}}\end{bmatrix},$
(14)
where the $k\times k$ matrix $\Phi_{\mathrm{DC}}$ and the $k\times(d-k)$
matrix $\Delta_{\mathrm{DC}}$ are drawn from the sub-matrices of $\Phi$ and
$\Delta$. Then the (13) can be rewritten as
$C_{\mathrm{DC}}=\begin{bmatrix}\Phi_{\mathrm{DC}}S+\Delta_{\mathrm{DC}}T^{t}&\Phi_{\mathrm{DC}}T\end{bmatrix}.$
(15)
The $C_{\mathrm{DC}}$ is split into two parts
$C_{\mathrm{DC}}=[\begin{array}[]{cc}C_{\mathrm{DC}}^{\Phi}&C_{\mathrm{DC}}^{\Delta}\end{array}]$,
where the $k$-column part $C_{\mathrm{DC}}^{\Phi}$ corresponds to
$\Phi_{\mathrm{DC}}S+\Delta_{\mathrm{DC}}T^{t}$, and the $(d-k)$-column part
$C_{\mathrm{DC}}^{\Delta}$ corresponds to $\Phi_{\mathrm{DC}}T$. Then the (15)
is reformulated as
$C_{\mathrm{DC}}^{\Phi}=\Phi_{\mathrm{DC}}S+\Delta_{\mathrm{DC}}T^{t};$ (16)
$C_{\mathrm{DC}}^{\Delta}=\Phi_{\mathrm{DC}}T.$ (17)
As the $\Phi_{\mathrm{DC}}$ is non-singular by the second condition of the
encoding matrix, the DC can compute the matrix
$T=\Phi_{\mathrm{DC}}^{-1}C_{\mathrm{DC}}^{\Delta}$, and subsequently, the
$S=\Phi_{\mathrm{DC}}^{-1}(C_{\mathrm{DC}}^{\Phi}-\Delta_{\mathrm{DC}}T^{t})$.
#### 2.2.3 Partial downloading scheme
Chen and Wang [18] indicate that the above data reconstruction process
involves a certain amount of redundancy. In the data reconstruction process,
the DC completely downloads $k$ vectors
$\\{\underline{c}_{i_{j}}^{t}|j=1,\ldots,k\\}$ with length $d$ for each
$\underline{c}_{i_{j}}^{t}$, to be used to reconstruct the $B={k+1\choose
2}+k(d-k)$ message symbols. As $(kd-B)={k\choose 2}\geq 0$, this process
potentially downloads ${k\choose 2}$ redundant symbols. To avoid the wasted
transmission resource, the [18] develops a partial downloading scheme on the
Exact-MBR code. By the scheme, the DC can only download the
$C_{\mathrm{DC}}^{\Delta}$ and the upper triangular part of
$C_{\mathrm{DC}}^{\Phi}$. Totally, the DC exactly download $B$ symbols.
In data reconstruction process, the sub-matrix $T$ can be solved by the
equality (17). Let
$D_{\mathrm{DC}}=C_{\mathrm{DC}}^{\Phi}-\Delta_{\mathrm{DC}}T^{t}$ (18)
denote the solvable part in (16). Thus, the (16) is rewritten as
$\Phi_{\mathrm{DC}}S=D_{\mathrm{DC}}.$ (19)
In the scheme [18], the DC only downloads the upper triangular part of
$C_{\mathrm{DC}}^{\Phi}$, so the upper triangular part of $D_{\mathrm{DC}}$ is
also accessible. The main idea of solving (19) is to utilize the symmetry of
$S$. The process can be divided into $k$ stages, and each stage solves a
column of $S$ in the backward order. While the $d$-th column of $S$ have been
solved, the $d$-th row of $S$ is also obtained by symmetry of $S$. The
obtained $d$-th row of $S$ will be utilized in the later decoding stages. By
such recursive decoding process, a symmetric matrix $S$ can be completely
solved.
## 3 Repair-by-transfer codes
This section proposes a new class of $(n,k,d=n-1)$ repair-by-transfer codes at
(5). Upon describing the code constructions, two basic entities, termed as the
message matrix $\hat{M}$ and the encoding matrix $\hat{\Phi}$, are defined as
follows. The $\hat{M}$ is a $n\times n$ matrix constructed from two sub-
matrices $\hat{S}$ and $\hat{T}$. The $\hat{S}$ is a $k\times k$ skew-
symmetric matrix determined by ${k\choose 2}$ message symbols. The skew-
symmetric matrix is defined as a square matrix $A$ satisfying $A=-A^{t}$. For
each entry $A[i,j]$ in the skew-symmetric matrix, the equality holds
$A[i,j]=-A[j,i]$. Note that the diagonal entries of skew-symmetric matrix $A$
are filled with zeros $A[i,i]=0$. By the above definition, the strictly upper
triangular part of $\hat{S}$ (excluding the diagonal entries) is filled with
${k\choose 2}$ message symbols, and the lower triangular part $\hat{S}$ is
filled with the corresponding values such that the skew symmetric condition
holds. The remaining $B-{k\choose 2}=k(n-k)$ message symbols are formed as the
second matrix $\hat{T}$ with $k\times(n-k)$. The $n\times n$ message matrix
$\hat{M}$ is defined as
$\hat{M}=\begin{bmatrix}\hat{S}&\hat{T}\\\
-\hat{T}^{t}&\mathbf{0}\end{bmatrix},$ (20)
where the $\mathbf{0}$ denotes a $(n-k)\times(n-k)$ zero matrix. Notably, the
$\hat{M}$ is also a skew-symmetric matrix.
For the encoding matrix, this matrix is defined as a $n\times n$ square matrix
of the form
$\hat{\Psi}=\begin{bmatrix}\hat{\Phi}&\hat{\Delta}\end{bmatrix},$ (21)
where the size of the matrix $\hat{\Phi}$ is $n\times k$, and the size of
matrix $\hat{\Delta}$ is $n\times(n-k)$. The $\hat{\Psi}$ is chosen in such a
way that
i) Any $k$ rows of $\hat{\Phi}$ are linearly independent;
ii) The matrix $\hat{\Psi}$ is non-singular.
The above conditions can be met by choosing $\hat{\Phi}$ to be a $n\times k$
Vandermonde matrix, and the $\hat{\Delta}$ is defined as
$\hat{\Delta}=\begin{bmatrix}\mathbf{0}\\\ I_{n-k}\end{bmatrix},$ (22)
where the $\mathbf{0}$ is a $k\times(n-k)$ zero matrix, and the $I_{n-k}$ is a
$(n-k)\times(n-k)$ identity matrix. By above definitions, the feasible range
of $n$ is $n\leq q$ over $GF(q)$. Furthermore, the $\hat{\Phi}$ can adopt the
extended Vandermonde matrix, which is the encoding matrix of the $(q+1,k)$
doubly extended RS code, as the form. Then the $n$ can be extended to $n=q+1$.
By above matrices, the construction of repair-by-transform code is formulated
as a congruence
$\hat{C}=\hat{\Psi}\hat{M}\hat{\Psi}^{t}.$
There is a useful theorem used in the code constructions: The $n\times n$
matrix $\hat{C}$ congruent to a skew-symmetric matrix $\hat{M}$ is also skew-
symmetric. Next, we modify the $\hat{C}$ to obtain a symmetric one
$\check{C}$. Each entry in strictly lower triangular part of $\hat{C}$ is
replaced with its negation value, resulting in a symmetric matrix $\check{C}$.
Equivalently, for each row $\hat{c}_{j}^{t}$ in $\hat{C}$, a modified row
$\check{c}_{j}^{t}$ in $\check{C}$ is obtained by assigning each entry to
$\check{c}_{j}^{t}[i]=\begin{cases}\hat{c}_{j}^{t}[i]&\textrm{if }i\geq j;\\\
-\hat{c}_{j}^{t}[i]&\textrm{otherwise.}\end{cases}$ (23)
The output $\check{C}$ is the generated codewords. The $n$ rows of $\check{C}$
are then respectively stored in $n$ distinct nodes. For $1\leq i\leq n$, the
$n$-element row $\check{c}_{i}^{t}$ is stored in an individual network node
indexed as $i$. As the diagonal entries
$\\{\check{c}_{i}^{t}[i]=0\\}_{i=1}^{n}$ are always zeros, those zero symbols
do not require storage space. Thus, each node takes $n-1$ units of memory
space to store a row of $\check{C}$, and the parameter configuration (5) holds
$\alpha=n-1$.
_Example 1:_ We give an example for $(n=5,k=3)$ repair-by-transfer codes over
$GF(4)$. By (5), other parameters are set as $d=\alpha=4$, $\beta=1$, and
$B=9$. By the definition of message matrix (20), the matrix $\hat{M}$ is
filled with 5 message symbols $\\{u_{i}\\}_{i=1}^{9}$ as follows:
$\hat{M}=\begin{bmatrix}0&u_{1}&u_{2}&u_{3}&u_{4}\\\
-u_{1}&0&u_{5}&u_{6}&u_{7}\\\ -u_{2}&-u_{5}&0&u_{8}&u_{9}\\\
-u_{3}&-u_{6}&-u_{8}&0&0\\\ -u_{4}&-u_{7}&-u_{9}&0&0\end{bmatrix}.$
As $-u_{i}=u_{i}$ over the field of characteristic two, the $\hat{M}$ is also
a symmetric matrix. For the encoding matrix, the matrix $\hat{\Phi}$ is chosen
as the $5\times 3$ extended Vandermonde matrix given by
$\hat{\Phi}=\begin{bmatrix}1&0&0\\\ 1&1&1\\\ 1&\omega&\omega^{2}\\\
1&\omega^{2}&\omega^{4}\\\ 0&0&1\end{bmatrix},$
where the $\omega$ denotes the primitive element of $GF(4)$. By the
$\hat{\Delta}$ defined in (22), the encoding matrix is expressed as
$\hat{\Psi}=\begin{bmatrix}1&0&0&0&0\\\ 1&1&1&0&0\\\
1&\omega&\omega^{2}&0&0\\\ 1&\omega^{2}&\omega^{4}&1&0\\\
0&0&1&0&1\end{bmatrix}.$
As the $\hat{M}$ is skew-symmetric, the congruence
$\hat{C}=\hat{\Psi}\hat{M}\hat{\Psi}^{t}$ is also skew-symmetric, expressed as
$\hat{C}=\begin{bmatrix}0&c_{1}&c_{2}&c_{3}&c_{4}\\\
-c_{1}&0&c_{5}&c_{6}&c_{7}\\\ -c_{2}&-c_{5}&0&c_{8}&c_{9}\\\
-c_{3}&-c_{6}&-c_{8}&0&c_{10}\\\ -c_{4}&-c_{7}&-c_{9}&-c_{10}&0\end{bmatrix}.$
(24)
Then each entry of strictly lower triangular part of $\hat{C}$ is replaced
with its additive inverse value, resulting in
$\check{C}=\begin{bmatrix}0&c_{1}&c_{2}&c_{3}&c_{4}\\\
c_{1}&0&c_{5}&c_{6}&c_{7}\\\ c_{2}&c_{5}&0&c_{8}&c_{9}\\\
c_{3}&c_{6}&c_{8}&0&c_{10}\\\ c_{4}&c_{7}&c_{9}&c_{10}&0\end{bmatrix}.$ (25)
Notably, as the $\hat{C}$ is over the field of characteristic two, the
$\check{C}=\hat{C}$ can be directly obtained without any arithmetic
operations.
### 3.1 Node-repairing process
The node-repairing process utilizes the symmetry of $\check{C}$. Suppose the
node $h_{0}$ fails, and the failure node stores the vector
$\check{c}_{h_{0}}^{t}$ at the $h_{0}$-th row of $\check{C}$. By the symmetry
of $\check{C}$, the $h_{0}$-th row of $\check{C}$ is equivalent to the
$h_{0}$-th column of $\check{C}$, whose entries (excluding the entry at main
diagonal) are respectively stored in $n-1$ non-failure nodes. Thus, the
replacement node can directly download the elements at the $h_{0}$-th column
$\check{C}$ from other $n-1$ nodes. Let the $\check{c}_{j}^{t}[i]$ denote the
$i$-th element of the row $\check{c}_{j}^{t}$. The formulation is given by
$\check{c}_{h_{0}}^{t}[i]=\begin{cases}0&\textrm{if }i=h_{0};\\\
\check{c}_{i}^{t}[h_{0}]&\textrm{Otherwise.}\end{cases}$ (26)
Consequently, this node-repairing process does not involve any arithmetic
operations at the helper nodes and the replacement node, as illustrated in
Example 1. In (25), if any one row of $\check{C}$ is erased, this row can be
regenerated through the aid of corresponding column in $\check{C}$.
### 3.2 Data reconstruction process with full downloading
In data reconstruction, the DC accesses the $k$ rows
$\\{\check{c}_{i_{1}}^{t},\check{c}_{i_{2}}^{t},\ldots,\check{c}_{i_{k}}^{t}\\}$,
which are respectively downloaded from $k$ connected nodes
$\\{i_{1},i_{2},\ldots,i_{k}\\}$. To begin with, each row
$\check{c}_{i_{j}}^{t}$, $1\leq j\leq k$, is restored to the original vector
$\hat{c}_{i_{j}}^{t}$ via the inversion of formula (23). The restored results
$\\{\hat{c}_{i_{1}}^{t},\ldots,\hat{c}_{i_{k}}^{t}\\}$ are formed as a
$k\times n$ matrix $\hat{C}_{\mathrm{DC}}$ which is a sub-matrix of $\hat{C}$.
By construction, the $\hat{C}_{\mathrm{DC}}$ possesses the equality
$\hat{C}_{\mathrm{DC}}=\hat{\Psi}_{\mathrm{DC}}\hat{M}\hat{\Psi}^{t},$ (27)
where the $k\times(n-1)$ matrix
$\hat{\Psi}_{\mathrm{DC}}=\begin{bmatrix}\hat{\Phi}_{\mathrm{DC}}&\hat{\Delta}_{\mathrm{DC}}\end{bmatrix}=\left[\begin{array}[]{c}\psi_{i_{1}}^{t}\\\
\vdots\\\ \psi_{i_{k}}^{t}\end{array}\right]$
consists of the $k$ encoding rows of $\hat{C}_{\mathrm{DC}}$.
As the $\hat{\Psi}$ is non-singular by the second condition of $\hat{\Psi}$,
the $\hat{C}_{\mathrm{DC}}$ in (27) is then post-multiplied by its inversion
$(\hat{\Psi}^{t})^{-1}$, to obtain a $k\times n$ matrix
$\hat{D}_{\mathrm{DC}}=\hat{C}_{\mathrm{DC}}(\hat{\Psi}^{t})^{-1}=\hat{\Psi}_{\mathrm{DC}}\hat{M}.$
(28)
The term $\hat{\Psi}_{\mathrm{DC}}\hat{M}$ in (28) is then decomposed as two
parts:
$\hat{\Psi}_{\mathrm{DC}}\hat{M}=\begin{bmatrix}\hat{\Phi}_{\mathrm{DC}}\hat{S}-\hat{\Delta}_{\mathrm{DC}}\hat{T}^{t}&\hat{\Phi}_{\mathrm{DC}}\hat{T}\end{bmatrix}.$
To elaborate the process, the $\hat{D}_{\mathrm{DC}}$ is split into two parts
$\hat{D}_{\mathrm{DC}}=[\begin{array}[]{cc}\hat{D}_{\mathrm{DC}}^{\Phi}&\hat{D}_{\mathrm{DC}}^{\Delta}\end{array}]$,
where the left part $\hat{D}_{\mathrm{DC}}^{\Phi}$ has $k$ columns and the
right part $\hat{D}_{\mathrm{DC}}^{\Delta}$ has $(n-k)$ columns, so
$\hat{D}_{\mathrm{DC}}^{\Phi}=\hat{\Phi}_{\mathrm{DC}}\hat{S}-\hat{\Delta}_{\mathrm{DC}}\hat{T}^{t};$
(29) $\hat{D}_{\mathrm{DC}}^{\Delta}=\hat{\Phi}_{\mathrm{DC}}\hat{T}.$ (30)
By the first definition of $\hat{\Psi}$, the $\hat{\Phi}_{\mathrm{DC}}$ is
non-singular. Thus, the DC can compute the matrix
$\hat{T}=\hat{\Phi}_{\mathrm{DC}}^{-1}\hat{D}_{\mathrm{DC}}^{\Delta}$; and
subsequently, the
$\hat{S}=\hat{\Phi}_{\mathrm{DC}}^{-1}(\hat{D}_{\mathrm{DC}}^{\Phi}+\hat{\Delta}_{\mathrm{DC}}\hat{T}^{t})$.
### 3.3 Systematic version of repair-by-transfer codes
To construct a systematic version of repair-by-transfer codes, a message-
symbol remapping procedure is employed to determine the entries of $\hat{M}$.
Without loss of generality, we declare that the source data are embedded in
the first $k$ rows of $\hat{C}$. To reduce the computational cost, the
$\hat{\Phi}$, which is the sub-matrix of $\hat{\Psi}$, is defined as the
encoding matrix of $(n,k)$ systematic RS codes. The matrix contains two parts
expressed as
$\hat{\Phi}=\begin{bmatrix}I_{k}\\\ \check{\Phi}\end{bmatrix},$ (31)
where the first $k$ rows of $\hat{\Phi}$ is an identity matrix $I_{k}$. By the
$\hat{\Delta}$ defined in (22), the encoding matrix $\hat{\Psi}$ is thus
formulated as
$\hat{\Psi}=\begin{bmatrix}I_{k}&\mathbf{0}\\\
\check{\Phi}&I_{n-k}\end{bmatrix}.$
By the above encoding matrix, the encoding formula
$\hat{C}=\hat{\Psi}\hat{M}\hat{\Psi}^{t}$ can be rewritten as
$\hat{C}=\begin{bmatrix}S&S\check{\Phi}^{t}+T\\\
\check{\Phi}S-T^{t}&\check{\Phi}S\check{\Phi}^{t}+\check{\Phi}T-T^{t}\check{\Phi}^{t}\end{bmatrix}.$
To achieve the systematic condition, the first $k$ rows of $\hat{C}$,
expressed as $\begin{bmatrix}S&S\check{\Phi}^{t}+T\end{bmatrix}$, are defined
as the source data. Let
$U=[\begin{array}[]{cc}U_{\mathrm{L}}&U_{\mathrm{R}}\end{array}]$ denote a
$k\times n$ matrix consisting of $B$ source symbols. The $U_{\mathrm{L}}$ is a
$k\times k$ skew-symmetric matrix whose strictly upper-triangular part is
filled with ${k\choose 2}$ source symbols, and other entries are filled with
the corresponding values to satisfy the skew-symmetry condition. The
$U_{\mathrm{R}}$ is a $k\times(n-k)$ matrix filled with $k\times(n-k)$ source
symbols. The systematic condition gives two equations
$S=U_{\mathrm{L}};S\check{\Phi}^{t}+T=U_{\mathrm{R}}.$
By above two equations, the $\hat{C}$ can be rewritten as
$\hat{C}=\begin{bmatrix}U_{\mathrm{L}}&U_{\mathrm{R}}\\\
-U_{\mathrm{R}}^{t}&V\end{bmatrix},$
where the $V$ is a $(n-k)\times(n-k)$ matrix defined as
$V=\check{\Phi}U_{\mathrm{R}}-U_{\mathrm{R}}^{t}\check{\Phi}^{t}-\check{\Phi}U_{\mathrm{L}}\check{\Phi}^{t}$.
As other three parts of $\hat{C}_{\mathrm{DC}}$, namely $U_{\mathrm{L}}$ and
$\pm U_{\mathrm{R}}$, are available without the arithmetic computations, the
matrix $V$ is the remaining unknown objective to be computed. It is noted that
the matrix $V$ is a skew-symmetric matrix, so is the $\hat{C}$. The
computation of $T$ involves the matrix product $\check{\Phi}U_{\mathrm{R}}$
and the congruence $\check{\Phi}U_{\mathrm{L}}\check{\Phi}^{t}$, and the term
$U_{\mathrm{R}}^{t}\check{\Phi}^{t}$ can be directly obtained via transposing
the result $\check{\Phi}U_{\mathrm{R}}$. As the $\check{\Phi}$ identifies the
encoding matrix of parity part in the $(n,k)$ systematic RS codes, the product
$\check{\Phi}U_{\mathrm{R}}$ denotes the parity parts of RS codes for each
column of $U_{\mathrm{R}}$. For the congruence
$\check{\Phi}U_{\mathrm{L}}\check{\Phi}^{t}$, the $(n,k)$ systematic RS
encoding is applied on each column of $U_{\mathrm{L}}$ to obtain the parity
part $\check{\Phi}U_{\mathrm{L}}$. Then the $(n,k)$ systematic RS encoding is
applied on each row of $\check{\Phi}U_{\mathrm{L}}$, resulting in the
$\check{\Phi}U_{\mathrm{L}}\check{\Phi}^{t}$ at the parity part. By above
steps, the product $\check{\Phi}U_{\mathrm{R}}$ requires $O(k(n-k)^{2})$
operations, and the transformation
$\check{\Phi}U_{\mathrm{L}}\check{\Phi}^{t}$ requires $O(2k^{2}(n-k))$
operations.
### 3.4 Partial downloading scheme
For the data reconstruction in Sec. 3.2, we suppose that the DC completely
downloads the $k$ vectors $\\{\check{c}_{i_{j}}^{t}|j=1,\ldots,k\\}$, and the
length of each vector $\check{c}_{i_{j}}^{t}$ is $n-1$. Thus, the total number
of downloaded symbols is $(n-1)k$, which is much larger than the size of
message $B=(n-1)k-{k\choose 2}$. By utilizing the symmetry of $\check{C}$, the
DC can exactly download $B$ symbols.
For any two distinct codeword vectors $\check{c}_{i_{j}}^{t}$ and
$\check{c}_{i_{l}}^{t}$ in $\check{C}$, we have
$\check{c}_{i_{j}}^{t}[i_{l}]=\check{c}_{i_{l}}^{t}[i_{j}]$ by the symmetric
property, so the DC can download this symbol only from either the node $i_{j}$
or the node $i_{l}$. Based on this observation, the $k$ connected nodes can
avoid the total of ${k\choose 2}$ symbols to be transmitted. An simple
transmission strategy is that, the first node $i_{1}$ transmits the whole
$n-1$ symbols $\check{c}_{i_{1}}^{t}$ to the DC. Then the second node $i_{2}$
can only transmit $n-2$ symbols of $\check{c}_{i_{2}}^{t}$ to the DC, as the
symbol $\check{c}_{i_{2}}^{t}[i_{1}]=\check{c}_{i_{1}}^{t}[i_{2}]$ does not
need to be transmitted. Inductively, the connected node $i_{j}$ can only
transmit $n-j$ symbols of $\check{c}_{i_{j}}^{t}$ to the DC, for
$j=1,\ldots,n$. The above policy is simple, but the data throughputs for each
node is imbalanced. Thus, an alternative transmission policy is presented in
the following. It is noted that each node can save $(k-1)/2$ symbols of data
transmission on average, and this value is achieved for odd $k$ by the
proposed transmission policy. For even $k$, as the value $(k-1)/2$ is not an
integer, the proposed transmission policy can save $k/2-1$ symbols in each
odd-index node, and $k/2$ symbols in each even-index node.
Given any two connected nodes with indices $i_{j}$, $i_{l}$ and $1\leq i,l\leq
k$, we define a decision criterion as
$D(j,l)=\begin{cases}min\\{j,l\\}&\textrm{if }j+l\textrm{ is even};\\\
max\\{j,l\\}&\textrm{otherwise.}\end{cases}$ (32)
As any two distinct nodes $i_{j}$ and $i_{l}$ simultaneously store a common
symbol, the $D(j,l)\in\\{j,l\\}$ returns the index of the chosen node to avoid
the transmission of this common symbol. Hence the DC downloads this element
from another un-chosen node. Two examples are given in Figure 2 tabulating the
exhaustive outputs of $D(i,j)$ for $k=5$ and $6$. In the case $k=5$, each node
omits two symbols in transmission. In the case $k=6$, the nodes
$\\{g_{1},g_{3},g_{5}\\}$ omit two symbols in transmission, and the nodes
$\\{g_{2},g_{4},g_{6}\\}$ omit three symbols in transmission.
The valid of decision criterion (32) is explained as follows. Given a node
indexed by $X$, we consider the output of $D(X,y)$ for $y=1,\ldots,k$. If the
(32) outputs $X=D(X,y)$ for a specific $y$, the node $X$ can omit the
transmission of a symbol, and DC will download this symbol from another node
$y$. To satisfy the equality $X=D(X,y)$, the range of $y$ are drawn from
$y\in\\{\ldots,X-1-2i,\ldots,X-1,X+2,\ldots,X+2i,\ldots\\}$ and $1\leq y\leq
k$. Thus, there are about $k/2$ distinct symbols of $y$, and the condition for
bandwidth balance holds.
| 5 | 4 | 3 | 2
---|---|---|---|---
1 | 1 | 4 | 1 | 2
2 | 5 | 2 | 3 |
3 | 3 | 4 | |
4 | 5 | | |
(a)
| 6 | 5 | 4 | 3 | 2
---|---|---|---|---|---
1 | 6 | 1 | 4 | 1 | 2
2 | 2 | 5 | 2 | 3 |
3 | 6 | 3 | 4 | |
4 | 4 | 5 | | |
5 | 6 | | | |
(b)
Figure 2: Two examples of the outputs of decision criterion $D(j,l)$. (a)
$k=5$. (b) $k=6$.
## 4 Systematic Exact-MBR coding algorithm
Based on the framework of $(n,k,d)$ Exact-MBR codes [9] in Sec. 2.2, this
section presents a systematic form of encoding matrix $\Psi$, where the
feasible range of $n$ are $n\leq q$ over $GF(q)$. Then the partial downloading
scheme is developed on the proposed Exact-MBR codes. Upon describing the
proposed encoding matrix, the encoding (6) can be divided into $\alpha$
individual columns given by
$\overline{c}_{i}=\Psi\overline{m}_{i},$ (33)
where the $\overline{m}_{i}$ indicates the $i$-th column of $M$, and the
result $\overline{c}_{i}$ is the $i$-th column in $C$. The (33) can be
rewritten as
$\overline{c}_{i}=\Psi\overline{m}_{i}=\begin{bmatrix}\Phi&\Delta\end{bmatrix}\left[\begin{array}[]{c}\overline{m}_{i}^{a}\\\
\overline{m}_{i}^{b}\end{array}\right]=\Phi\overline{m}_{i}^{a}+\Delta\overline{m}_{i}^{b},$
(34)
where the $\overline{m}_{i}^{a}$ denotes the $k$-element vector located in the
upper part of the $\overline{m}_{i}$, and the $\overline{m}_{i}^{b}$ denotes
the remaining $(d-k)$-elements located in the lower part of the
$\overline{m}_{i}$.
By the first condition of Exact-MBR encoding matrix, the $\overline{m}_{i}$
can be reconstructed from arbitrary $d$ elements in $\overline{c}_{i}$. By the
second condition, if the term $\Delta\overline{m}_{i}^{b}$ is given, the
$\overline{m}_{i}^{a}$ can be reconstructed from arbitrary $k$ elements in
$\overline{c}_{i}$. Under above observations, Section 4.1 presents a class of
modified version of Reed-Solomon codes, termed as partially systematic Reed-
Solomon (PSRS) codes, to satisfy those conditions. Section 4.2 shows that the
encoding matrix of the systematic Exact-MBR codes. Section 4.3 presents the
partial downloading scheme.
### 4.1 Partially systematic Reed-Solomon codes
We define the partially systematic Reed-Solomon (PSRS) code associated with
three parameters $(n,k,d)$ where $k\leq d<n$. The $n$ is the codeword length,
the $d$ is the message length, and the $k$ is the length of systematic part.
The input is expressed as a $d$-element vector
$\underline{c}=\begin{bmatrix}\underline{a}&\underline{b}\end{bmatrix}$, where
the sub-vector $\underline{a}=[a_{1}...a_{k}]$ denotes the $k$ systematic
symbols, and the sub-vector $\underline{b}=[b_{1}...b_{d-k}]$ denotes the
remaining $d-k$ non-systematic symbols. By definition, the systematic part
$\underline{a}$ is embedded in the first $k$ elements of the generated
codeword. This subsection presents the constructions of $(n,k,d)$ PSRS codes
via the polynomial evaluation approach. Let the $G(x)$ denote the coding
polynomial constructed from the message $\underline{c}$. The degree of $G(x)$
is $deg(C(x))<d$. The codeword symbols are the evaluations of $C(x)$ at $n$
distinct points:
$\\{C(x_{1}),C(x_{2}),\ldots,C(x_{n})\\}.$ (35)
As the code is over $GF(q)$, the code suffices for $n\leq q$. By the partially
systematic condition, the first $k$ codeword symbols are equivalent to the
systematic message symbols. Thus,
$C(x_{i})=a_{i},\forall i=1,2,\ldots,k.$ (36)
In the following, the $C(x)$ is properly defined to satisfy the partial
systematic condition.
The $C(x)$ is defined as the sum of two polynomials
$C(x)=\Phi(x)+\Delta(x),$ (37)
where the polynomial $\Phi(x)$ is constructed from $\underline{a}$, and the
$\Delta(x)$ is constructed from $\underline{b}$. The $\Phi(x)$, and
$deg(\Phi(x))<k$, is defined as
$\Phi(x)=\sum_{i=1}^{k}a_{i}\prod_{j\neq i}\frac{x-x_{j}}{x_{i}-x_{j}}.$ (38)
This follows the form of Lagrange polynomial. Thus, the $\Phi(x)$ possesses
the systematic property:
$\Phi(x_{i})=a_{i},\forall i=1,2,\ldots,k.$
The polynomial $\Delta(x)$ is defined as the multiplication of two
polynomials:
$\Delta(x)=\Gamma(x)B(x).$ (39)
The polynomial $\Gamma(x)$ has $k$ roots located in the evaluation points of
systematic part:
$\Gamma(x)=\prod_{i=1}^{k}(x-x_{i}).$ (40)
The $B(x)$ is constructed from the $(d-k)$-element vector $\underline{b}$. The
$B(x)$ can be chosen as the systematic or non-systematic form. For example, a
non-systematic form with geometric progression is expressed as
$B(x)=\sum_{i=1}^{d-k}b_{i}x^{i-1}.$ (41)
By the above definitions, it can be shown that the partial systematic
condition (36) holds:
$\displaystyle C(x_{i})$ $\displaystyle=\Phi(x_{i})+\Delta(x_{i})$ (42)
$\displaystyle=\Phi(x_{i})+0\times B(x_{i})=a_{i},\forall i=1,2,\ldots,k.$
In summary, the encoding algorithm includes four major steps listed as
follows:
i). Compute the coefficients of $\Phi(x)$.
ii). Compute the product $\Delta(x)=\Gamma(x)B(x)$, where the coefficients of
$\Gamma(x)$ can be computed in advance.
iii). Compute the summation $C(x)=\Phi(x)+\Delta(x)$.
iv). Evaluate the values $\\{C(x_{1}),\ldots,C(x_{n})\\}$ to obtain the
codeword symbols.
If the encoding algorithm is implemented in the native way, the computational
complexities of the four steps are $O(k^{2})$, $O(k(d-k))$, $O(k)$, and
$O(dn)$, respectively.
To reduce the complexity complexity, we observe that the fast Fourier
transforms (FFT) can be utilized to reduce the computational cost in steps
(i), (ii) and (iv). The conceptual ideas are addressed below. In step (i), the
(38) can be calculated via fast Lagrange interpolation [22] with complexity
$O(k\log^{2}k)$. Alternatively, the fast Reed-Solomon encoding algorithms can
also be used in (38). If the code is operated on Fermat field
$GF(q+1),q\in\\{2,4,16,65536\\}$, the (38) can be calculated via inverse fast
Fourier transform with complexity $O(k\log k)$ (see [24] and [25]). If the
code is operated on finite field with characteristic two
$GF(q),q\in\\{2,4,8,\ldots\\}$, the [23] proposed an coding algorithm with
complexity $O(q\log^{2}q)$. The step (ii) is a polynomial multiplication. By
using FFT, the complexity can be reduced to $O(d\log d)$. In step (iv), the
polynomial evaluations can be computed with FFT, and the complexity is
$O(n\log n)$.
It is noted that the PSRS codes can also be implemented with generator
polynomials. The details are placed in appendix.
#### 4.1.1 Full erasure decoding from $d$ codeword symbols
The message vector $\underline{c}$ can be reconstructed from arbitrary $d$ out
of $n$ codeword symbols $\\{y_{i}=C(z_{i})|1\leq i\leq d\\}$. By the subset of
codeword symbols, the $C(x)$ is constructed via Lagrange interpolation:
$C(x)=\sum_{i=1}^{d}y_{i}\prod_{j\neq i}\frac{x-z_{j}}{z_{i}-z_{j}}.$ (43)
The $C(x)$ is then divided by $\Gamma(x)$ to obtain a quotient $B(x)$ and a
remainder $\Phi(x)$. The $k$ evaluations $a_{i}=\Phi(x_{i})$, $1\leq i\leq k$,
are the systematic part $\underline{a}$, and the coefficients of $B(x)$ are
the non-systematic part $\underline{b}$.
#### 4.1.2 Partial erasure decoding from $k$ codeword symbols
Suppose the non-systematic part $\underline{b}$ is given. In this case, we
shows that the systematic part $\underline{a}$ can be reconstructed from
arbitrary $k$ out of $n$ codeword symbols $\\{y_{i}=C(z_{i})|1\leq i\leq
k\\}$. By the given $\underline{b}$, the polynomial $\Delta(x)$ can be
constructed. Then the $k$ evaluation values of $\Phi(x)$ are calculated via
$\Phi(z_{i})=C(z_{i})-\Delta(z_{i}),\forall i=1...k.$ (44)
By the $k$ evaluation values of $\Phi(x)$, the $\Phi(x)$ can be interpolated
via Lagrange polynomial, and the $\underline{a}$ is the $k$ evaluations
$a_{i}=\Phi(x_{i})$.
### 4.2 Encoding matrix of proposed Exact-MBR codes
As the $(n,k,d)$ PSRS codes satisfy the conditions of Exact-MBR codes, the
encoding matrix of $(n,k,d)$ PSRS codes can be chosen as the $\Psi$. For the
systematic part $\underline{a}$, the coding polynomial $\Phi(x)$ formulates a
generator matrix corresponding to the component $\Phi$ in encoding matrix
$\Psi$. By the definition of $\Phi(x)$, the entries of matrix $\Phi$ are
$\Phi[l,i]=\prod_{j=1;j\neq i}^{k}\frac{x_{l}-x_{j}}{x_{i}-x_{j}},\textrm{for
}i=1,\ldots,k.$ (45)
Consequently, the first $k$ rows of $\Phi$ is a $k\times k$ identity matrix
$I_{k}$. For the non-systematic part $\underline{b}$, the coding polynomial
$\Delta(x)$ formulates a generator matrix corresponding to the component
$\Delta$ in encoding matrix $\Psi$. By the definition of $\Delta(x)$, the
entries of matrix $\Delta$ are
$\Delta[l,i]=x_{l}^{(i-1)}\Gamma(x_{l})\textrm{, for }i=1,\ldots,d-k.$ (46)
As $\Gamma(x_{l})=0$ for $1\leq l\leq k$, the first $k$ rows of $\Delta$ are
entirely filled with zeros. Then the encoding matrix $\Psi$ is obtained by
combining the $\Phi$ and $\Delta$. Thus, the first $k$ rows of $\Psi$ are in
the form $\begin{bmatrix}I_{k}&\mathbf{0}\end{bmatrix}$, so that the
corresponding first $k$ rows of the code matrix $C$ are expressed as
$[\begin{array}[]{cc}S&T\end{array}]$. Hence, the proposed Exact-MBR code is
systematic. As stated previously, the proposed $\Psi$ satisfies the two
conditions of Exact-MBR encoding matrix, which enables the node-repairing
algorithm and data reconstruction algorithm addressed in Sec. 2.2.
_Example 2:_ We give an example for $(n=6,k=3,d=4)$ Exact-MBR codes over
$GF(7)$. By (4), other parameters are set as $\alpha=4$, $\beta=1$, and $B=9$.
By the definition of message matrix (8), the matrices $M$ is filled with 9
message symbols $\\{u_{i}\\}_{i=1}^{9}$. The $S$, $T$ and $M$ are given by
$\displaystyle S$ $\displaystyle=\begin{bmatrix}u_{1}&u_{2}&u_{3}\\\
u_{2}&u_{5}&u_{6}\\\ u_{3}&u_{6}&u_{8}\end{bmatrix},T=\begin{bmatrix}u_{4}\\\
u_{7}\\\ u_{9}\end{bmatrix};$ $\displaystyle M$
$\displaystyle=\begin{bmatrix}u_{1}&u_{2}&u_{3}&u_{4}\\\
u_{2}&u_{5}&u_{6}&u_{7}\\\ u_{3}&u_{6}&u_{8}&u_{9}\\\
u_{4}&u_{7}&u_{9}&0\end{bmatrix}.$
The coding polynomial $C(x)$ of $(n=6,k=3,d=4)$ PSRS code is chosen as
$\displaystyle\Phi(x)=$ $\displaystyle
a_{1}\times\frac{(x-2)(x-3)}{2}+a_{2}\times\frac{(x-1)(x-3)}{6}$
$\displaystyle+a_{3}\times\frac{(x-1)(x-2)}{2};$ $\displaystyle\Delta(x)=$
$\displaystyle(x-1)(x-2)(x-3)b_{1}.$
By above definitions, the corresponding matrices $\Phi$ and $\Delta$ are as
follows:
$\Phi=\begin{bmatrix}1&0&0\\\ 0&1&0\\\ 0&0&1\\\ 1&4&3\\\ 3&6&6\\\
6&6&3\end{bmatrix};\Delta=\begin{bmatrix}0\\\ 0\\\ 0\\\ 6\\\ 3\\\
4\end{bmatrix}.$
The encoding matrix $\Psi=\begin{bmatrix}\Phi&\Delta\end{bmatrix}$ is the
combination of $\Phi$ and $\Delta$.
### 4.3 Partial downloading scheme
This subsection presents the partial downloading scheme on the proposed
systematic Exact-MBR codes. Similar to the [18], the proposed scheme only
downloads the entire $C_{\mathrm{DC}}^{\Delta}$ and the lower (or upper,
alternatively) triangular part of $C_{\mathrm{DC}}^{\Phi}$. Precisely, each
connected node $i_{j}$ passes a portion of the code fragment $c_{i_{j}}^{t}$
in the lower/upper triangular part of $C_{\mathrm{DC}}$. By (17), the $T$ can
be successfully solved. Then the lower/upper triangular part of
$D_{\mathrm{DC}}$ can be computed via (18). The two cases are respectively
considered as follows.
#### 4.3.1 Data collector downloads the lower triangular part of
$C_{\mathrm{DC}}^{\Phi}$
In this case, the DC can access the lower triangular part of
$D_{\mathrm{DC}}$. The computational structure can be divided into $k$ stages,
and the $l$-th stage solves the $l$-th column $\overline{s}_{l}$ of $S$. In
the first stage, as the first column of $D_{\mathrm{DC}}$ are fully located in
the lower triangular part of $D_{\mathrm{DC}}$, the first column
$\overline{s}_{1}$ of $S$ can be solved successfully. By the symmetry of $S$,
the first row of $S$ is also obtained
$\underline{s}_{1}=\overline{s}_{1}^{t}$. Let $\underline{i}_{l}^{t}$ denote a
row vector with one at the $l$-th position and zeros elsewhere. By the
definition of proposed encoding matrix, the obtained $\underline{s}_{1}$ is at
the first row (systematic part) of $\Phi$. Thus, we have the equation
$\underline{i}_{1}^{t}S=\underline{s}_{1}$ which will be utilized in the
upcoming decoding stages.
In the $l$-th stage, $1\leq l\leq k$, the DC can access the
$\\{\underline{d}_{i_{l}}^{t}[l],\ldots,\underline{d}_{i_{k}}^{t}[l]\\}$ in
the $l$-th column of lower triangular part of $D_{\mathrm{DC}}$, and the
corresponding encoding rows are
$\\{\underline{\phi}_{i_{l}}^{t},\ldots,\underline{\phi}_{i_{k}}^{t}\\}$. In
the previous stages, we obtain $l-1$ equations:
$\underline{i}_{j}^{t}\overline{s}_{l}=\underline{s}_{j}[l],\forall
j=1,\ldots,l-1.$
It is noted that the $\\{\underline{i}_{j}^{t}|1\leq j\leq l-1\\}$ are the
first $l-1$ rows of $\Phi$. The above equations are combined to obtain
$\begin{bmatrix}\underline{i}_{1}^{t}\\\ \vdots\\\ \underline{i}_{l-1}^{t}\\\
\underline{\phi}_{i_{l}}^{t}\\\ \vdots\\\
\underline{\phi}_{i_{k}}^{t}\end{bmatrix}\overline{s}_{l}=\begin{bmatrix}\underline{s}_{1}[l]\\\
\vdots\\\ \underline{s}_{l-1}[l]\\\ \underline{d}_{i_{l}}^{t}[l]\\\ \vdots\\\
\underline{d}_{i_{k}}^{t}[l]\end{bmatrix}.$ (47)
Let the $D_{1}^{l}$ denote the matrix at the left-hand-side of (47). To solve
the $\overline{s}_{l}$ successfully, the $D_{1}^{l}$ should be non-singular.
Then we have $\underline{s}_{l}=\overline{s}_{l}^{t}$, and the
$\underline{i}_{l}^{t}S=\underline{s}_{l}$ can be utilized in the upcoming
decoding stages.
The non-singularity of $D_{1}^{l}$ is discussed below. In the $D_{1}^{l}$, the
set
$\underline{i}^{t}=\\{\underline{i}_{1}^{t},\ldots,\underline{i}_{l-1}^{t}\\}$
are the first $l-1$ rows of $\Phi$, and the set
$\underline{\phi}^{t}=\\{\underline{\phi}_{i_{l}}^{t},\ldots,\underline{\phi}_{i_{k}}^{t}\\}$
are $k-l$ rows in $\Phi$. As any $k$ rows of $\Phi$ are non-singular, the
$D_{1}^{l}$ is also non-singular, as long as the two sets are mutually
exclusive $\underline{i}^{t}\cap\underline{\phi}^{t}=\emptyset$. To satisfy
this condition, the order of fragments $[g_{1},\ldots,g_{k}]$ in
$C_{\mathrm{DC}}$ should follow a special condition: For the systematic
fragment $\underline{c}_{l}^{t}$, $1\leq l\leq k$, downloaded from the node
$i_{j}$, the $\underline{c}_{l}^{t}$ is placed at the $g_{j}$-th row of
$C_{\mathrm{DC}}$, where $g_{j}\leq l$.
_Example 3:_ By following the codes given by Example 1, we assume that the DC
connects to nodes 1, 2, and 4 respectively corresponding to encoding rows
$\begin{bmatrix}1&0&0&0\end{bmatrix}$, $\begin{bmatrix}0&1&0&0\end{bmatrix}$
and $\begin{bmatrix}1&4&3&6\end{bmatrix}$. The three rows of $C_{\mathrm{DC}}$
are arranged as
$C_{\mathrm{DC}}=\begin{bmatrix}1&0&0&0\\\ 0&1&0&0\\\ 1&4&3&6\end{bmatrix}M.$
The DC downloads the whole $C_{\mathrm{DC}}^{\Delta}$ and the lower triangular
part of $C_{\mathrm{DC}}^{\Phi}$. The $C_{\mathrm{DC}}^{\Delta}$ possesses the
equation given by
$C_{\mathrm{DC}}^{\Delta}=\begin{bmatrix}1&0&0\\\ 0&1&0\\\
1&4&3\end{bmatrix}T.$
By the equation, the $T$ can be solved to obtain
$\\{\tilde{u}_{4},\tilde{u}_{7},\tilde{u}_{9}\\}$, where the tilde symbol
$\tilde{\bullet}$ indicates the solved terms. By the solved $T$, the DC
calculates the lower triangular part of $D_{\mathrm{DC}}$ via
$D_{\mathrm{DC}}=C_{\mathrm{DC}}^{\Phi}-\begin{bmatrix}0\\\ 0\\\
6\end{bmatrix}\begin{bmatrix}\tilde{u}_{4}&\tilde{u}_{7}&\tilde{u}_{9}\end{bmatrix}.$
Let $D[i,j]$ denote the entry of $D_{\mathrm{DC}}$ at the $i$-th row and
$j$-th column. The accessible part of $D_{\mathrm{DC}}$ is
$\begin{bmatrix}D[1,1]&-&-\\\ D[2,1]&D[2,2]&-\\\
D[3,1]&D[3,2]&D[3,3]\end{bmatrix}=\begin{bmatrix}1&0&0\\\ 0&1&0\\\
1&4&3\end{bmatrix}S,$
where the notation ”$-$” indicates the inaccessible entries. Firstly, by the
first column of $D_{\mathrm{DC}}$, the first column of $S$ is solved. The
solved symbols possesses the equality:
$\begin{bmatrix}\tilde{u}_{1}&\tilde{u}_{2}&\tilde{u}_{3}\end{bmatrix}=\begin{bmatrix}1&0&0\end{bmatrix}S.$
(48)
Secondly, to decode the second column of $S$, we have
$\begin{bmatrix}\tilde{u}_{2}\\\ D[2,2]\\\
D[3,2]\end{bmatrix}=\begin{bmatrix}1&0&0\\\ 0&1&0\\\
1&4&3\end{bmatrix}\begin{bmatrix}u_{2}\\\ u_{5}\\\ u_{6}\end{bmatrix}.$
Then the symbols $\\{u_{5},u_{6}\\}$ are solved. The solved symbols possess
the equality:
$\begin{bmatrix}\tilde{u}_{2}&\tilde{u}_{5}&\tilde{u}_{6}\end{bmatrix}=\begin{bmatrix}0&1&0\end{bmatrix}S.$
(49)
By the third column of $D_{\mathrm{DC}}$ and the (48)(49), we have
$\begin{bmatrix}\tilde{u}_{3}\\\ \tilde{u}_{6}\\\
D[3,3]\end{bmatrix}=\begin{bmatrix}1&0&0\\\ 0&1&0\\\
1&4&3\end{bmatrix}\begin{bmatrix}u_{3}\\\ u_{6}\\\ u_{8}\end{bmatrix}.$
Then the symbol $u_{8}$ is solved successfully.
#### 4.3.2 Data collector downloads the upper triangular part of
$C_{\mathrm{DC}}^{\Phi}$
In this case, the DC accesses the upper triangular part of $D_{\mathrm{DC}}$
defined in (19). The steps are very similar to the above decoding scheme. The
decoding structure can be expressed as $k$ stages, and each stage extracts a
column of $S$ in backward order. That is, the $l$-th stage extracts the
$(k+1-l)$-th column $\overline{s}_{k+1-l}$ of $S$. In the $l$-th stage, the DC
can access the
$\\{\underline{d}_{i_{1}}^{t}[l],\ldots,\underline{d}_{i_{k+1-l}}^{t}[l]\\}$
taken from the $(k+1-l)$-th column of $D_{\mathrm{DC}}$ in upper triangular
part, and the corresponding encoding rows are
$\\{\underline{\phi}_{i_{1}}^{t},\ldots,\underline{\phi}_{i_{k+1-l}}^{t}\\}$.
Furthermore, we also have $l-1$ equations by the previous stages:
$\underline{i}_{j}^{t}\overline{s}_{l}=\underline{s}_{j}[l],\forall
j=k+2-l,\ldots,k.$
Those equations are combined to obtain
$\begin{bmatrix}\underline{\phi}_{i_{1}}^{t}\\\ \vdots\\\
\underline{\phi}_{i_{k+1-l}}^{t}\\\ \underline{i}_{k+2-l}^{t}\\\ \vdots\\\
\underline{i}_{k}^{t}\end{bmatrix}\overline{s}_{k+1-l}=\begin{bmatrix}\underline{d}_{i_{1}}^{t}[l]\\\
\vdots\\\ \underline{d}_{i_{k+1-l}}^{t}[l]\\\ \underline{s}_{k+2-l}[l]\\\
\vdots\\\ \underline{s}_{k}[l]\end{bmatrix}.$ (50)
Let the $D_{2}^{l}$ denote the left-hand-side matrix in (50). To decode the
$\overline{s}_{k+1-l}$, the $D_{2}^{l}$ should be non-singular, and this
condition induces that
$\\{\underline{\phi}_{i_{1}}^{t},\ldots,\underline{\phi}_{i_{k+1-l}}^{t}\\}\cap\\{\underline{i}_{k+2-l}^{t},\ldots,\underline{i}_{k}^{t}\\}=\emptyset$,
for $1\leq l\leq k$. By the above condition, the systematic fragment
$\underline{c}_{l}^{t}$ downloaded from the node $i_{j}$ is placed at the
$g_{j}$-th row of $C_{\mathrm{DC}}$, where $1\leq l\leq g_{j}\leq k$. Then the
$\overline{s}_{k+1-l}$ can be solved successfully, and the formula
$\underline{i}_{k+1-l}^{t}S=\overline{s}_{k+1-l}^{t}$ is utilized in the
upcoming decoding stages.
#### 4.3.3 The time-sharing policy to balance the bandwidth requirements on
each connected node
In the above two partial downloading schemes, both partial downloading schemes
have the disadvantage that the transmission amounts for $k$ connected nodes
are excessively unbalanced. To overcome this drawback, we can iteratively
switch the two partial downloading schemes during the whole transmission
rounds. Specifically, if a node $i_{j}$ transmits the elements of a code
fragment in the lower triangular of $C_{\mathrm{DC}}^{\Phi}$ at this
transmission round, this node will transmit the elements of next code fragment
in the upper triangular of $C_{\mathrm{DC}}^{\Phi}$ at the next transmission
round. By this time-sharing policy, each node transmits $d-(k-1)/2$ symbols in
each transmission round on average.
As stated in Sections 4.3.1 and 4.3.2, the two partial downloading schemes
respectively give two different conditions on the order $[g_{1},\ldots,g_{k}]$
of the downloaded fragments in $C_{\mathrm{DC}}$. Since the time-sharing
policy iteratively applies two partial downloading schemes, the two conditions
should be satisfied simultaneously. The intersection of two conditions is
that, the systematic fragment $\underline{c}_{l}^{t}$ downloaded from the node
$i_{j}$ is placed at the $g_{j}$-th row of $C_{\mathrm{DC}}$, where $1\leq
g_{j}=l\leq k$.
## 5 Comparisons and discussions
In this section, we compare the proposed codes with prior works. The results
are briefly summarized in Tables I and II.
### 5.1 Comparisons for Repair-by-transfer codes
TABLE I: Comparisons for repair-by-transfer codes over $GF(q)$. | Down. policy | Range of $n$ | Enc. comp.
---|---|---|---
Shah et al. [10] | - | ${n\choose 2}\leq q+1$ | $O(n^{4})$
Ours (Section 3) | Partial | $n\leq q+1$ | $O(n^{3})$
Figure 3: The simulations of [10] and the proposed repair-by-transfer codes.
This subsection compares the proposed repair-by-transfer codes with the [10]
introduced in Sec. 2.1. As shown in Sec. 2.1, the field size of is at least
${n\choose 2}\leq N+1$. For the proposed repair-by-transfer codes, Section 3.3
states that the feasible range of $n$ can be extended up to $n\leq N+1$ via
the extended Vandermonde matrix. Hence we conclude that the size of finite
field is significantly reduced.
Another issue is the computational complexities. We compare the complexities
of both codes over the same finite field $GF(q)$. For the [10], it is evident
that the $(N,K)=\left({n\choose 2},B\right)$ MDS code dominates the whole
computational overhead. By employing the $\left({n\choose 2},B\right)$
systematic RS code, the encoding complexity is given by $O(({n\choose
2}-B)B)$. For the proposed scheme, the systematic version Sec. 3.3 computes
the matrix $V$, whose computational cost is dominated by two terms
$\check{\Phi}U_{\mathrm{R}}$ and the
$\check{\Phi}U_{\mathrm{L}}\check{\Phi}^{t}$. As stated in Sec. 3.3, both
terms take a total of $O(k(n-k)^{2})+O(2k^{2}(n-k))=O(k(n^{2}-k^{2}))$
operations. To magnify the difference between both codes further, we consider
the case $k=cn$ with a constant $c$. In this case, the big-O representation of
both codes are simplified into $O(({n\choose 2}-B)B)=O(n^{4})$ , and
$O(k(n^{2}-k^{2}))=O(n^{3})$, respectively. Thus, the proposed code reduces
one order of magnitude in big-O complexity representation. The real
simulations of two codes are shown in Fig. 3. Both codes are written in JAVA,
and the programs are running on Intel i7-950, 4GB RAM, Windows 8. We test the
case $k=n/2$ at $n=\\{8,12,\ldots,64\\}$ over $GF(2^{16})$. In the simulation,
the source data are generated by a random number generator. The Y-axis
represents the logarithm of the encoding time of the $B$ input symbols on
average. As shown in Fig. 3, the performance of the proposed codes is better
than the [10] if the $n$ is larger than a specific value. For the small value
of $n$, we conjecture that the structure of [10] is more simple, and the
proposed algorithm contains a number of redundant arithmetic operations in the
the congruence $V$, so that the [10] is better.
### 5.2 Comparisons for systematic Exact-MBR codes
TABLE II: Comparisons for Exact-MBR codes over $GF(q)$. | Syst. | Down. policy | Range of $n$ | Enc. complexity
---|---|---|---|---
Rashmi et al. [9] | N | Full | $n\leq q$ | $O(nd^{2})$
Y | Full | $n\leq(q+k-d)$ or $n\leq q$ | $O(nd^{2})$
Gong and Wang[18] | N | Partial | $n\leq q$ | $O(nd^{2})$
Ours (Section 4) | Y | Partial | $n\leq q$ | $O(nd^{2})$ or $O(n\log n)$
In the following, we compare the proposed systematic Exact-MBR codes with the
[9], in terms of the range of $n$ and the encoding complexity. For the range
of $n$, the [9] presents two distinct forms for the encoding matrix, so the
$n$ has two distinct upper bounds. The first form is expressed as
$\Psi=\left[\begin{array}[]{cc}I_{k}&\mathbf{0}\\\
\tilde{\Phi}&\tilde{\Delta}\end{array}\right],$ (51)
where $I_{k}$ denotes a $k\times k$ identity matrix, $\mathbf{0}$ is a
$k\times(d-k)$ zero matrix. The
$[\begin{array}[]{cc}\tilde{\Phi}&\tilde{\Delta}\end{array}]$ is a
$(n-k)\times d$ Cauchy matrix, where the sizes of $\tilde{\Phi}$ and
$\tilde{\Delta}$ are $(n-k)\times k$ and $(n-k)\times(d-k)$, respectively. As
stated by [9], the (51) meets the two conditions of Exact-MBR encoding matrix.
By definition, a $(n-k)\times d$ Cauchy matrix requires $n-k+d$ distinct
symbols. As the $GF(q)$ contains a total of $q$ distinct symbols, the feasible
range of $n$ is
$n-k+d\leq q\Rightarrow n\leq q+k-d.$
As addressed in Sec. 4.1, the range of $n$ for the PSRS codes is $n\leq q$, so
is the proposed Exact-MBR code. Due to $k\leq d$, the proposed codes have
larger range of $n$.
In the second form of encoding matrix [9], the range of $n$ is also $n\leq q$.
However, the second form is not explicit and the matrix generation requires an
additional matrix inversion and multiplication step. An explicit form can
facilitate the further development on the codes. For example, the partial
decoding algorithm proposed in Sec. 4.3 is based on the observations on the
form of encoding matrix. If the encoding matrix is not explicit, the partial
decoding algorithm may become more difficult to be designed. Furthermore, by
Appendix, the proposed $(n,k)$ PSRS codes can be implemented by generator
polynomials. The size of generator polynomial is $(n-k)$, which is lower than
the size of encoding matrix $(n-k)\times k$ in parity part. Thus, the
generator polynomial approach is more common in usage.
The encoding complexities of those codes are discussed below. Suppose those
three codes are implemented with native matrix product approach. As the sizes
of encoding matrix and message matrix are $n\times d$ and $d\times d$ for the
three codes, the encoding complexity is $O(nd^{2})$. Furthermore, Section 4.1
indicates that the PSRS codes can be implemented with fast Fourier transforms.
By FFTs, the encoding complexity can be reduced to $O(n\log n)$.
### 5.3 Comparisons for partial downloading schemes on Exact-MBR codes
The partial downloading scheme is useful to reduce the requisite throughput to
reconstruct data. This subsection highlights the differences between [18] and
ours. First, the proposed scheme requires that the systematic fragments should
be placed at a specific row of $C_{\mathrm{DC}}$. On the other hand, the [18]
do not require this condition as the [18] is developed on non-systematic
codes. Second, in our survey, this is the first work of considering the
throughput balance on the connected nodes.
### 5.4 Simulations for systematic Exact-MBR codes
Figure 4: The simulations of systematic Exact-MBR codes with native approach
and fast approach.
As stated in Section 4.1, the PSRS codes can be implemented with FFT. By
employing the fast algorithm of PSRS code, we expected that the encoding time
of Exact-MBR codes can be reduced. Based on this motivation, we implement the
native and fast approaches of Exact-MBR codes, and the simulation results are
shown in Figure 4. Both codes are written in JAVA, and the programs are
running on Intel i7-950, 4GB RAM, Windows 8. We test the case $k=3/8\times n$,
$d=n/2$, at $n\in\\{32,64,128,256,512\\}$ over Fermat field $GF(2^{16}+1)$. As
shown in Figure 4, the fast approach works better for larger $n$. Otherwise,
the native approach is suggested.
## 6 Conclusions
The contributions of this paper can be organized in two parts. First, a new
class of repair-by-transfer codes are proposed at $d=n-1$ MBR points. As
compared with prior works, the proposed repair-by-transfer code demands
smaller finite field and lower big-O complexity. The partial downloading
scheme is also developed on the proposed repair-by-transfer codes to avoid the
unnecessary symbol transmissions. The simulation shows that the proposed
repair-by-transfer codes require fewer arithmetic operations than the prior
work when $n$ is larger than a specific value. Second, for all feasible
parameters $(n,k,d)$, we present an encoding matrix for systematic Exact-MBR
codes via the partially systematic Reed-Solomon codes. To minimize the number
of transmitted symbols in data reconstruction process, the partial downloading
scheme is designed on the proposed Exact-MBR codes. However, the transmission
amount for those connected nodes are excessively unbalanced. Thus, a time-
sharing scheme is presented to balance the bandwidths requirements on those
connected nodes. The proposed Exact-MBR codes can be implemented via fast
Fourier transforms. As shown in the simulations, the fast approach has better
encoding performance for large $n$.
## Appendix A Partially systematic Reed-Solomon codes by generator polynomial
The appendix presents another approach of $(n,k,d)$ PSRS codes by generator
polynomials. In this approach, the messages and codewords are formulated as
polynomials. Thus, the message $\underline{a}$ and $\underline{b}$ are
$a(x)=a_{0}+a_{1}x+\ldots+a_{k-1}x^{k-1};$
$b(x)=b_{0}+b_{1}x+\ldots+b_{d-k-1}x^{d-k-1}.$
The codeword polynomial is defined as
$c(x)=c_{0}(x)+c_{1}(x),$
where the $c_{0}(x)$ is the codeword generated from $a(x)$ via $(n,k)$
systematic RS code, and the $c_{1}(x)$ is the codeword generated from $b(x)$
via $(n-k,d-k)$ RS code. Precisely, for the construction of $c_{0}(x)$, the
generator polynomial of $(n,k)$ systematic RS code is defined as
$g_{0}(x)=(x-1)(x-\alpha)\ldots(x-\alpha^{n-k-1}).$
Then the parity polynomial $r_{0}(x)$ is calculated through polynomial
division
$r_{0}(x)=x^{n-k}a(x)\textrm{ (mod }g_{0}(x)).$
The codeword $c_{0}(x)$ is expressed as the concatenation of $a(x)$ and
$r_{0}(x)$:
$c_{0}(x)=x^{n-k}a(x)-r_{0}(x).$ (52)
For the construction of $c_{1}(x)$, the generator polynomial of $(n-k,d-k)$ RS
code is defined as
$g_{1}(x)=(x-1)(x-\alpha)\ldots(x-\alpha^{n-d-1}).$
The $c_{1}(x)$ can be formed as the systematic or non-systematic version. For
the systematic case, the codeword polynomial is defined as
$r_{1}(x)=x^{n-d}b(x)\textrm{ (mod }g_{1}(x));$
$c_{1}(x)=x^{n-d}b(x)-r_{1}(x).$ (53)
The polynomial $a(x)$ is embedded in the $c(x)$ between $x^{n-k}$ and
$x^{n-1}$, as the degree of $c_{1}(x)$ is less than $n-k$. Thus, the partially
systematic condition holds. By generator polynomial, the length of this
$(n,k,d)$ coding algorithm gets up to $n\leq q-1$ over $GF(q)$. The decoding
algorithms are explained in the following.
### A.1 Full erasure decoding from $d$ codeword symbols
The $a(x)$ and $b(x)$ can be reconstructed by arbitrary $d$ out of $n$
coefficients of the $c(x)$. As $g_{0}(x)$ and $g_{1}(x)$ are respectively the
factors of $c_{0}(x)$ and $c_{1}(x)$, the $\gcd(g_{0}(x),g_{1}(x))=g_{0}(x)$
is also the factor of the $c(x)$. Therefore the $(n,k,d)$ PSRS code is
isomorphic to the $(n,d)$ RS code with the generator polynomial $g_{0}(x)$.
Thus, the $c(x)$ can be reconstructed from arbitrary $d$ out of $n$
coefficients via Forney algorithm. Forney algorithm is a method to compute the
erasures of BCH codes at known error locations. When the $c(x)$ is completely
reconstructed, the $a(x)$ is located in the systematic part of $c(x)$. Then
the $c_{0}(x)$ can be computed from $a(x)$, and subsequently the
$c_{1}(x)=c(x)-c_{0}(x)$. Thus, the $b(x)$ is decoded from $c_{1}(x)$.
### A.2 Partial erasure decoding from $k$ codeword symbols
Given the $b(x)$, the message $a(x)$ can be reconstructed by arbitrary $k$ out
of $n$ coefficients in $c(x)$. By (53), the $c_{1}(x)$ is calculated from
$b(x)$. Since we have $k$ coefficients in $c(x)$, the corresponding $k$
coefficients in $c_{0}(x)=c(x)-c_{1}(x)$ can also be calculated. As the
$c_{0}(x)$ is the codeword of $(n,k)$ systematic RS code, the $c_{0}(x)$ can
be completely recovered via Forney algorithm. Then the message $a(x)$ is
obtained from $c_{0}(x)$.
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* [25] S. J. Lin and W. H. Chung, ”An Efficient (n, k) Information Dispersal Algorithm based on Fermat Number Transforms,” to appear in IEEE Transactions on Information Forensics and Security, doi: 10.1109/TIFS.2013.2270892.
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* [27] O. Khan, R. Burns, J. S. Plank, W. Pierce and C. Huang, ”Rethinking Erasure Codes for Cloud File Systems: Minimizing I/O for Recovery and Degraded Reads,” FAST 2012: 10th USENIX Conference on File and Storage Technologies, San Jose, CA, Feb. 2012.
| Sian-Jheng Lin was born in Taichung, Taiwan, in 1981. He received the B.S.,
M.S., and Ph.D. degrees in computer science from National Chiao Tung
University, in 2004, 2006, and 2010, respectively. He is currently a
postdoctoral fellow with the Research Center for Information Technology
Innovation, Academia Sinica. His recent research interests include data hiding
and error control coding.
---|---
| Wei-Ho Chung was born in Kaohsiung, Taiwan, in 1978. He received the B.Sc.
and M.Sc. degrees in Electrical Engineering from National Taiwan University,
Taipei City, Taiwan, in 2000 and 2002 respectively. From 2005 to 2009, he was
with the Electrical Engineering Department at University of California, Los
Angeles, where he obtained his Ph.D. degree. From 2000 to 2002, he worked on
routing protocols in the mobile ad hoc networks in the M.Sc. program in
National Taiwan University. From 2002 to 2005, he was a system engineer at
ChungHwa Telecommunications Company, where he worked on data networks. In
2008, he was an research intern working on CDMA systems in Qualcomm, Inc. From
2007 to 2009, he was a Teaching Assistant at UCLA. From June to December 2009,
Dr. Chung had been working as a research associate in San Diego, California,
on wireless communications for multimedia communications and unequal error
protection for video transmission. His research interests include
communications, signal processing, and networks. Dr. Chung received the Taiwan
Merit Scholarship from 2005 to 2009, and the Best Paper Award in IEEE WCNC
2012. Dr. Chung has been an assistant research fellow in Research Center for
Information Technology Innovation in Academia Sinica, Taiwan, since January
2010.
---|---
|
arxiv-papers
| 2013-09-15T11:54:39 |
2024-09-04T02:49:50.927875
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Sian-Jheng Lin and Wei-Ho Chung",
"submitter": "Sian-Jheng Lin",
"url": "https://arxiv.org/abs/1309.3752"
}
|
1309.3838
|
# Linear isometries between real JB∗-triples and C∗-algebras
Maria Apazoglou Queen Mary, University of London, London, United Kingdom
[email protected] and Antonio M. Peralta Departamento de Análisis
Matemático, Universidad de Granada,
Facultad de Ciencias 18071, Granada, Spain [email protected]
###### Abstract.
Let $T:A\to B$ be a (not necessarily surjective) linear isometry between two
real JB∗-triples. Then for each $a\in A$ there exists a tripotent $u_{a}$ in
the bidual, $B^{\prime\prime},$ of $B$ such that
1. $(a)$
$\\{u_{a},T(\\{f,g,h\\}),u_{a}\\}=\\{u_{a},\\{T(f),T(g),T(h)\\},u_{a}\\}$, for
all $f,g,h$ in the real JB∗-subtriple, $A_{a},$ generated by $a$;
2. $(b)$
The mapping $\\{u_{a},T(\cdot),u_{a}\\}:A_{a}\rightarrow B^{\prime\prime}$ is
a linear isometry.
Furthermore, when $B$ is a real C∗-algebra, the projection
$p=p_{a}=u_{a}^{*}u_{a}$ satisfies that $T(\cdot)p:A_{a}\rightarrow
B^{\prime\prime}$ is an isometric triple homomorphism. When $A$ and $B$ are
real C∗-algebras and $A$ is abelian of real type, then there exists a partial
isometry $u\in B^{\prime\prime}$ such that the mapping
$T(\cdot)u^{*}u:A\rightarrow B^{\prime\prime}$ is an isometric triple
homomorphism. These results generalise, to the real setting, some previous
contributions due to C.-H. Chu and N.-C. Wong, and C.-H. Chu and M. Mackey in
2004 and 2005. We give an example of a non-surjective real linear isometry
which cannot be complexified to a complex isometry, showing that the results
in the real setting can not be derived by a mere complexification argument.
Second author partially supported by the Spanish Ministry of Science and
Innovation, D.G.I. project no. MTM2011-23843, and Junta de Andalucía grants
FQM0199 and FQM3737.
Part of the results of this paper are contained in the first author’s PhD
thesis [1] presented to Queen Mary, University of London, funded by EPSRC. The
first author wishes to thank Professor Cho-Ho Chu for his valuable guidance
and support.
## 1\. Introduction
Despite many recent applications of real C∗-algebras in operator theory [8,
13], in JB∗-triples and infinite dimensional geometry [7, 8, 9, 17, 18, 22,
36, 37], the theory of real C∗-algebras has not been fully developed. In this
paper, we study the linear geometry of real C∗-algebras and in particular, the
not necessarily surjective linear isometries between real C∗-algebras, thereby
extending the isometry results in [10], [11] and [25] for complex C∗-algebras.
A celebrated result of Kadison [25] states that a surjective linear isometry
$T:A\rightarrow B$ between complex C∗-algebras is necessarily a Jordan triple
isomorphism, that is,
$T(ab^{*}c+cb^{*}a)=T(a)T(b)^{*}T(c)+T(c)T(b)^{*}T(a)\qquad(a,b,c\in A).$
This result need not be true when the hypothesis of $T$ being surjective is
not assumed. Nevertheless, it has been shown by C.-H. Chu and N.-C. Wong [9]
that even if $T$ is non-surjective, for each $a\in A$, there is a largest
projection $p_{a}\in B^{**}$, such that
$T(\cdot)p_{a}:A_{a}\rightarrow B^{**}$
is an isometric Jordan triple homomorphism, where $A_{a}$ denotes the
JB∗-subtriple of $A$ generated by $a$. In [10], C.-H. Chu and M. Mackey have
further extended this result to the case in which $A$ and $B$ are JB∗-triples,
showing, among other things, that, in this case, for each $a\in A$ there
exists a tripotent $u_{a}\in B^{**}$ such that
$\\{u_{a},T(\\{f,g,h\\}),u_{a}\\}=\\{u_{a},\\{T(f),T(g),T(h)\\},u_{a}\\}\ \
(f,g,h\in A_{a})$
and the mapping $\\{u_{a},T(\cdot),u_{a}\\}$ is an isometry.
Our objective is not only to extend these results to real C∗-algebras but also
to clarify the relationship between the projection $p_{a}$ and the tripotent
$u_{a}$ which was not given in [10]. Indeed, we show that both results above
are valid for real C∗-algebras and moreover, we have $p_{a}=u_{a}^{*}u_{a}$.
Prior to this result, we establish that when $A$ is a real C∗-algebra which
admits a complex character of real type, then for every not necessarily
surjective linear isometry $T$ from $A$ to another C∗-algebra $B$, there
exists a minimal projection $p$ in $B^{\prime\prime}$ satisfying:
1. $(a)$
$T\\{a,b,c\\}p=\\{T(a),T(b),T(c)\\}p$ and $pT(a)^{*}T(b)=T(a)^{*}T(b)p,$ for
all $a,b,c$ in $A$;
2. $(b)$
$T(\cdot)p:A\to B^{\prime\prime}$ is a non-zero triple homomorphism
(see Proposition 3.1 and Theorem 3.4). When $A$ does not admit complex
characters of real type, that is, when $A$ is of complex type we can find an
example of a linear surjection $T$ from $A$ into another real C∗-algebra such
that $T(\cdot)p=0$ whenever $T(\cdot)p$ is a triple homomorphism,
$T\\{a,b,c\\}p=\\{T(a),T(b),T(c)\\}p$ and $pT(a)^{*}T(b)=T(a)^{*}T(b)p,$ for
all $a,b,c$ in $A$ (compare Example 3.2).
We culminate the paper showing that every non-surjective linear isometry
between real JB∗-triples $T:A\to B$, although need not be a triple
homomorphism, locally reduces to a triple homomorphism, that is, for each
$a\in A$, there exists a tripotent $u\in B^{\prime\prime}$ such that
1. $(a)$
$\\{u,T(\\{f,g,h\\}),u\\}=\\{u,\\{T(f),T(g),T(h)\\},u\\}$, for all $f,g,h$ in
the real JB∗-subtriple generated by $a$;
2. $(b)$
The mapping $\\{u,T(\cdot),u\\}:A_{a}\rightarrow B^{\prime\prime}$ is a linear
isometry.
Furthermore, when $B$ is a real C∗-algebra, the projection $p=u^{*}u$
satisfies that $T(\cdot)p:A_{a}\rightarrow B^{\prime\prime}$ is an isometric
triple homomorphism (compare Theorem 3.9).
Our results cannot be achieved by simple complexification since, as shown
later, the complexification of a non-surjective real isometry need not be, in
general, an isometry. We develop some new techniques to accomplish our task.
## 2\. Preliminaries: Jordan structures in real C∗-algebras
Let $A$ be a real Banach*-algebra. Then $A$ is called a _real C ∗-algebra_ if
$\|a^{*}a\|=\|a\|^{2}$ and $1+a^{*}a$ is invertible for all $a\in A$. If $A$
is non-unital, then we require that $1+a^{*}a$ is invertible in the unit
extension $A\oplus\mathbb{R}1$. Equivalently, a real Banach *-algebra $A$ is a
real C∗-algebra if, and only if, it is isometrically *-isomorphic to a norm-
closed real *-algebra of bounded operators on a real Hilbert space (cf. [29,
Corollary 5.2.11]). Trivially, every (complex) C∗-algebra is a real C∗-algebra
when scalar multiplication is restricted to the real field. Let $A$ be a real
C∗-algebra. We denote by $A^{\prime}$ the dual space of $A$. The
complexification $A_{c}=A\oplus iA$ is a C∗-algebra in a suitable norm such
that $A$ identifies with a real closed *-subalgebra of $A_{c}$. The dual space
of a complex C∗-algebra $B$ will be denoted by $B^{*}$.
Throughout the paper, given a real or complex C∗-algebra $A$, $A_{sa}$ will
stand for the set of all self-adjoint elements in $A$.
We shall now survey some Jordan structures associated with real and complex
C∗-algebras. A (complex) _JB ∗-triple_ is a complex Banach space $A$ equipped
with a triple product $\\{\cdot,\cdot,\cdot\\}:A\times A\times A\rightarrow A$
which is linear and symmetric in the outer variables, conjugate linear in the
middle one and satisfies the following conditions:
1. (1)
(Jordan identity) for $a,b,x,y,z\in A$,
(2.1)
$\\{a,b,\\{x,y,z\\}\\}=\\{\\{a,b,x\\},y,z\\}-\\{x,\\{b,a,y\\},z\\}+\\{x,y,\\{a,b,z\\}\\};$
2. (2)
$D(a,a):A\rightarrow A$ is an hermitian linear operator with non-negative
spectrum, where $D(a,b)(x)=\\{a,b,x\\}$ with $a,b,x\in A$;
3. (3)
$\|\\{x,x,x\\}\|=\|x\|^{3}$ for all $x\in A$.
A complex subspace $B$ of $A$ is called a _subtriple_ if it is closed with
respect to the triple product, that is, $x,y,z\in B$ implies $\\{x,y,z\\}\in
B$. A real subtriple of $A$ is a real subspace of $A$ which is closed with
respect to the triple product. We define a _real JB ∗-triple_ to be a closed
real subtriple of a JB∗-triple. When restricted to real scalar multiplication,
every (complex) JB∗-triple is also a real JB∗-triple. It has been shown in
[22] that a real JB∗-triple can be complexified to become a JB∗-triple.
Furthermore, every real JB∗-triple $B$ is a real form of a complex JB∗-triple,
that is, there exist a (complex) JB∗-triple $B_{c}$ and a conjugate linear
isometry $\tau:B_{c}\rightarrow B_{c}$ of period 2 such that $B=\\{b\in
B_{c}\>:\>\tau(b)=b\\}$. We can actually identify $B_{c}$ with the
complexification of $B$.
A real C∗-algebra $A$ is a real JB∗-triple in the triple product:
(2.2) $\\{a,b,c\\}=\frac{1}{2}(ab^{*}c+cb^{*}a)\qquad(a,b,c\in A).$
We shall always assume the above triple product in a real or complex
C∗-algebra.
An element $e$ in a complex or real JB∗-triple $A$ is called _tripotent_ if
$\\{e,e,e\\}=e$. Let $A$ be a real JB∗-triple. For an element $a\in A$, we
define the operators $D(a):A\rightarrow A$ and $Q(a):A\rightarrow A$ by
$D(a)(x)=\\{a,a,x\\}\;\textrm{and}\;Q(a)(x)=\\{a,x,a\\}\qquad(x\in A).$
The Peirce projections associated with a tripotent $e\in A$ are given by
$\displaystyle P_{2}(e)$ $\displaystyle=Q(e)^{2}$ $\displaystyle P_{1}(e)$
$\displaystyle=2(D(e)-Q(e)^{2})$ $\displaystyle P_{0}(e)$
$\displaystyle=I-2D(e)+Q(e)^{2}.$
The space $A_{j}(e)=P_{j}(e)(A)$ is called the Peirce-$j$-space associated to
the tripotent $e$. Thus, each tripotent $e$ in a real or complex JB∗-tripe,
$A$, defines a decomposition of $A$ in terms of Peirce subspaces
$A=A_{0}(e)\oplus A_{1}(e)\oplus A_{2}(e).$
Another decomposition can be given in terms of the eigensubspaces of $Q(e)$,
$A=A^{1}(e)\oplus A^{-1}(e)\oplus A^{0}(e)$
where $A^{k}(e):=\\{x\in A:Q(e)(x):=\\{e,x,e\\}=kx\\}$ is a real Banach
subspace of $A$. Furthermore, the following identities and rules hold:
$\displaystyle A_{2}(e)=A^{1}(e)\oplus A^{-1}(e),$ $\displaystyle
A_{1}(e)\oplus A_{0}(e)=A^{0}(e)$
$\displaystyle\\{A^{i}(e),A^{j}(e),A^{k}(e)\\}\subseteq A^{ijk}(e),\hbox{
whenever }ijk\neq 0.$
A real JB∗-triple which is also a dual Banach space is called a real
(respectively, complex) _JBW ∗-triple_. Every real or complex JBW∗-triple has
a (unique) predual (cf. [3], [17, Corollary 3.2], [21] and [31]) and its
triple product is separately weak∗ continuous (see [31]).
Given a tripotent $e$ in a real or complex JB∗-triple $A$, the Peirce space
$A_{2}(e)$ is a real JB∗-algebra and $A^{1}(e)$ is a JB-algebra with product
and involution given by $a\circ_{e}b:=\\{a,e,b\\}$ and
$a^{\sharp_{e}}:=\\{e,a,e\\}$, respectively. It can be easily seen that, for
each $a\in A$, the expression
(2.3)
$P_{2}(e)\left\\{a,a,e\right\\}=\left\\{P_{2}(e)(a),P_{2}(e)(a),e\right\\}+\left\\{P_{1}(e)(a),P_{1}(e)(a),e\right\\}$
$=P^{1}(e)\left\\{a,a,e\right\\}$
defines a positive element in the JB-algebra $A^{1}(e)$ (see, for example the
arguments given in [5, Lemma 2.1], which remain valid in the real setting).
Two elements $a,b$ in a real or complex JB∗-triple $A$ are said to be
_orthogonal_ (written $a\perp b$) if $L(a,b)=0$. By Lemma 1 in [6] (whose
proof is also valid for real JB∗-triples), two elements $a,b$ in $A$ are
orthogonal if, and only if, $\left\\{a,a,b\right\\}=0$ (equivalently,
$\left\\{b,b,a\right\\}=0$). It follows from the Peirce rules that, for each
tripotent $e$ in $A$, the sets $A_{0}(e)$ and $A_{2}(e)$ are orthogonal.
In a real or complex JB∗-triple, $A$, a partial order relation, $\leq$, can be
defined in the set of all tripotents elements in $A$ given by $e\leq f$ if,
and only if, $f-e$ is a tripotent in $A$ with $f-e\perp e$. A non-zero
tripotent $e$ is called _minimal_ whenever $A^{1}(e)=\mathbb{R}e$ (since in
the complex case $A^{-1}(e)=iA^{1}(e)$, this definition is equivalent to say
$A_{2}(e)=\mathbb{C}e$). However, in the real setting the dimensions of
$A^{1}(e)$ and $A^{-1}(e)$ are not correlated as there exist examples of real
Cartan factors $A$ containing a minimal tripotent $e$ satisfying
dim$(A^{-1}(e))=\infty$. Another illustrating example is
$A=\mathbb{R},\mathbb{C}$ or $\mathbb{H}$ and $e=1$, then $e$ is minimal and
dim$(A^{-1}(e))=0,1$ and $3,$ respectively. In general, given a minimal
tripotent $e$ in a real JB∗-triple $A$, $\\{e,A,e\\}=A_{2}(e)$ need not
coincide with $\mathbb{R}e=A^{1}(e)$.
In [37] a tripotent $e$ is called minimal when it is minimal with respect to
the partial order relation $\leq$, that is, $e$ is minimal iff $e\geq f\neq 0$
implies $e=f$. To avoid confusion, let us call these kind of tripotents _order
minimal tripotents_. Clearly every minimal tripotent is order-minimal,
however, the reciprocal statement is not, in general, true. For example, the
unit element $1$ in $C[0,1]$ (real or complex valued) is order minimal because
there exists no tripotent $e\in C[0,1]$ satisfying $1>e$. When $A$ is a real
JBW∗-triple minimal and order minimal tripotents in $A$ coincide (cf.
Proposition 2.2 and Remark 2.3, [37]).
We recall that a minimal tripotent $e$ in a real JB∗-triple $E$ is called
_reduced_ whenever $E_{2}(e)=\mathbb{R}e$ (equivalently, $E^{-1}(e)=0$). The
real JB∗-triple $E$ is said to be _reduced_ when every minimal tripotent $e\in
E$ is reduced (cf. [30, 11.9]).
A real W∗-algebra is a real C∗-algebra $A$ having a predual,
$A_{{}_{{}^{\prime}}}$, in which case, $A_{{}_{{}^{\prime}}}$ is unique up to
linear isometric isomorphism and the product of $A$ is separately weak∗
continuous (cf. [24]). We note that the second dual $A^{\prime\prime}$ of a
real C∗-algebra $A$ is a real W∗-algebra. Moreover, let $A_{c}$ and
$A_{c}^{**}$ denote the complexifications of $A$ and its bidual, respectively,
then $A_{c}^{**}=A^{\prime\prime}\oplus iA^{\prime\prime}$. Every W∗-algebra
is a real JB∗-triple with respect to the product given in $(\ref{eq triple
product in C*-algebras})$.
An element $p$ in a real C∗-algebra $A$ is called a _projection_ if
$p=p^{*}=p^{2}$. A projection $p$ in $A$ is _minimal_ whenever it is a minimal
tripotent in $A$, that is, $A^{1}(p)=\\{x\in A:px^{*}p=x\\}=\mathbb{R}p$. It
is well known that, when $A$ is regarded as a real JB∗-triple, tripotents in
$A$ correspond to _partial isometries_ in $A$ (i.e. elements $u\in A$ with
$u^{*}u$ being a projection in $A$).
The next lemma should be known for experts but we were not able to find an
explicit reference.
###### Lemma 2.1.
Let $u$ be a tripotent in a real C∗-algebra $A$ regarded as a real JB∗-triple.
Then $u$ is minimal if, and only if, $p=u^{*}u$ (respectively, $q=uu^{*}$) is
a minimal projections in $A$.
###### Proof.
The equivalence follows from the following identities: $A^{1}(p)=pA_{sa}p,$
$uA^{1}(p)=A^{1}(u)=\\{x\in A:ux^{*}u=x\\}.$ ∎
Let $A$ be a real W∗-algebra. Two projections $p$ and $q$ in $A$ are called
_orthogonal_ if $pq=0$. A projection $p$ is _majorised_ by $q$, i.e. $p\leq q$
in the partial order of the cone $A_{+}$, if and only if $p(1-q)=0$. When $A$
is regarded as a JB∗-triple this partial order agrees with the partial order
defined in the set of tripotents (i.e. partial isometries) of $A$.
###### Lemma 2.2.
([29],Proposition 4.3.4) Let $A\subset B(H)$ be a real W∗-algebra. Then the
supremum $\bigvee_{i\in I}p_{i}$ of a family $\\{p_{i}\\}_{i\in I}$ of
projections in $A$ exists with respect to the ordering induced by the cone
$A_{+}$. Further, we have
$\bigvee_{i\in I}p_{i}:H\rightarrow\overline{[\cup_{i\in I}p_{i}H]}.$
is the projection onto the closed linear span $\overline{[\cup_{i\in
I}p_{i}H]}$ of $\cup_{i\in I}p_{i}H$.
###### Lemma 2.3.
Let $\\{p_{\phi}\\}_{\phi\in Q}$ be a family of projections in a real
W∗-algebra $W$. Let $p=\bigvee_{\phi\in Q}p_{\phi}$. If $ap_{\phi}=0$ for some
$a\in W$ and all $\phi\in Q$, then $ap=0$.
###### Proof.
By hypothesis, $a\in W$ annihilates on every element in the linear span of
$\bigcup_{i\in I}p_{i}(H)$, and by continuity $ap(\xi)=0$, for every $\xi\in
H$. ∎
If $B$ is a real JB∗-triple, then its second dual, $B^{\prime\prime}$, is also
a real JB∗-triple. Moreover, $B$ is the real form of a complex JB∗-triple
$B_{c}$, that is, there is a conjugate linear isometry $\tau:B_{c}\rightarrow
B_{c}$ of period 2 such that $B=\\{b\in B_{c}\>:\>\tau(b)=b\\}$. Further, the
bidual map $\sigma=\tau^{**}:B_{c}^{**}\rightarrow B_{c}^{**}$ is a conjugate
linear isometry of period 2 and $B^{\prime\prime}=\\{b\in
B_{c}^{**}\>:\>\sigma(b)=b\\}$. The “dual” map
$\widetilde{\tau}:B_{c}{{}^{*}}\rightarrow B_{c}^{*}$ defined by
$\widetilde{\tau}(\phi)(b)=\overline{\phi(\tau(b))}\qquad(\phi\in
B_{c}^{*},\,b\in B_{c})$
is a conjugate linear isometry of period 2 and the mapping
$(B_{c}^{*})^{\widetilde{\tau}}\to B^{\prime}$ $\varphi\mapsto\varphi|_{B}$
is a surjective linear isometry.
We observe that if $u$ is a tripotent in $B_{c}^{**}$, then $\sigma(u)$ is
also a tripotent in $B_{c}^{**}$. In fact, letting $U(B_{c}^{**})$ be the set
of all tripotents in $B_{c}^{**}$, the set $U(B^{\prime\prime})$ of tripotents
in $B^{\prime\prime}$ is $U(B_{c}^{**})\cap B^{\prime\prime}=\\{u\in
U(B_{c}^{**})\>:\>\sigma(u)=u\\}$ [17]. For a functional $\phi\in B_{c}^{*}$,
there exists a unique tripotent $u_{\phi}\in B_{c}^{**}$, called the _support
tripotent_ of $\phi$, such that $\phi=\phi\circ P_{2}(u_{\phi})$ and
$\phi|_{P_{2}(u_{\phi})(B_{c}^{**})}$ is a faithful normal positive functional
on the JBW∗-algebra $P_{2}(u_{\phi})(B_{c}^{**})$ (cf. [19]). We note that if
$u_{\phi}\in B_{c}^{**}$ is the support tripotent of $\phi\in B_{c}^{*}$, then
$\sigma(u_{\phi})$ is the support tripotent of $\sigma_{*}(\phi)$. On the
other hand, if $\phi$ is in $B^{\prime}$, its complex extension $\phi_{c}\in
B_{c}^{*}$ has support tripotent $u_{\phi_{c}}\in B_{c}^{**}$ such that
$\sigma(u_{\phi_{c}})=u_{\phi_{c}}$ since
$\sigma^{\sharp}(\phi_{c})=\phi_{c}$. Hence $u_{\phi_{c}}\in
U(B^{\prime\prime})$ and we call it the _support tripotent_ of $\phi$ in
$B^{\prime\prime}$, denoted by $u_{\phi}$ (cf. [36, Lemma 2.2]). Finally, we
note that $\phi$ is an extreme point of the closed unit ball of $A_{*}$ if and
only if its support tripotent $u_{\phi}$ is a minimal tripotent in $A$ [37,
Corollary 2.1]. Moreover, in this case $\phi=\phi P^{1}(u_{\phi})$ (and hence
$\phi=\phi Q(u_{\phi})$) ([37, Lemma 2.7]). Therefore
$\phi(x)u_{\phi}=\frac{1}{2}\left(Q(u_{\phi})^{2}(x)+Q(u_{\phi})(x)\right)=P^{1}(u_{\phi})(x),$
for every $x\in A$. Based on these observations and the complex results in [2,
Proposition 1.2] we have the following lemma.
###### Lemma 2.4.
([36, Lemma 2.4]) Let $B^{\prime\prime}$ be the second dual of a real
C∗-algebra $B$ and let $\phi\in B^{\prime}$ with $\|\phi\|=1$. Let $u,v\in
B^{\prime\prime}$ such that $\phi(u)=\|u\|=1=\phi(v)=\|v\|$ then
1. (1)
$\phi\\{x,y,v\\}=\phi\\{x,y,u\\}=\phi\\{y,x,u\\}$;
2. (2)
$\phi\\{x,x,u\\}\geq 0$;
3. (3)
$|\phi\\{x,y,u\\}|^{2}\leq\phi\\{x,x,u\\}\phi\\{y,y,u\\}$ (Cauchy-Schwarz
inequality),
for all $x,y\in B^{\prime\prime}$.
###### Lemma 2.5.
Let $B^{\prime\prime}$ be the second dual of a real C∗-algebra $B$ and let
$\phi\in B^{\prime}$ be an extreme point in the closed unit ball
$B^{\prime}_{1}$. Let $N_{\phi}=\\{b\in
B^{\prime\prime}\>:\>\phi\\{b,b,u\\}=0\\}$ where $u$ is the support tripotent
of $\phi$. Then $N_{\phi}=P_{0}(u)(B^{\prime\prime})$.
###### Proof.
By $(\ref{equation positive element})$ and the faithfulness of
$\phi|_{(B^{\prime\prime})^{1}(u)}$, we have $\\{b,b,u\\}=0$, that is, $b\perp
u$, and hence $b\in P_{0}(u)(B^{\prime\prime})$. ∎
A (real) linear functional $\varphi$ on a real C∗-algebra $A$ is said to be
positive when it is hermitian and maps positive elements in $A$ into
$\mathbb{R}^{+}_{0}$ (i.e. $\varphi(a)\geq 0$, for every positive element $a$
in $A$, and $\varphi(b^{*})=\varphi(b)$, for every $b$ in $A$). Following
standard notation, given a real C∗-algebra $A$, the quasi-state space of $A$,
$\mathcal{Q}(A),$ is defined as the set of all positive $\varphi\in
A^{\prime}$ satisfying $\|\varphi\|\leq 1$ (see, for example [29, §5.2]). It
is known that $\mathcal{Q}(A)$ is a weak∗ closed convex subset of
$(A^{\prime})_{1}$. We shall require later some results on the facial
structure of $\mathcal{Q}(A)$. Having in mind that $A_{sa}$ is a JC-algebra
with respect to the canonical Jordan product, the following Lemma was proved
in [32, Theorem 5.2].
Finally, let $A$ be a real or complex C∗-algebra or a JB-algebra. Following
standard notation (see [34, §3.11] and [14]), we recall that a projection $p$
in the bidual, $A^{**}$, of $A$ is said to be open if $A^{**}_{2}(p)\cap A$ is
weak∗-dense in $A^{**}_{2}(p)$. The projection $p$ is called closed when $1-p$
is open.
###### Lemma 2.6.
[32, Theorem 5.2] Let $\mathcal{Q}(A)$ be the quasi-state space of a real
C∗-algebra $A$. For every norm-closed face $F$ of $\mathcal{Q}(A)$ there
exists a (unique) closed projection $p\in A^{\prime\prime}$ satisfying that
$F=\\{\varphi\in\mathcal{Q}(A):\varphi(1-p)=0\\}.\hfill\Box$
## 3\. Isometries between real C∗-algebras
Our goal in this section is to show that any linear isometry $T:A\rightarrow
B$, surjective or not, between real C∗-algebras $A$ and $B$ reduces _locally_
to an (isometric) Jordan triple isomorphism by a projection $p=u^{*}u\in
B^{\prime\prime}$ for some partial isometry $u\in B^{\prime\prime}$.
In our first result we shall study contractive linear projections between real
C∗-algebras. In the complex setting, C.-H. Chu and N.-C. Wong proved in [11,
Proposition 2.2] that for every linear contraction $T$ between (complex)
C∗-algebras $A$ and $B$, there exists a largest projection $p\in B^{**}$ such
that
1. $(a)$
$T\\{a,a,a\\}p=\\{T(a),T(a),T(a)\\}p$;
2. $(b)$
$pT(a)^{*}T(a)=T(a)^{*}T(a)p$ for every $a\in A$.
A standard complex polarisation argument implies
1. $(a)$
$T\\{a,b,c\\}p=\\{T(a),T(b),T(c)\\}p$ and $pT(a)^{*}T(b)=T(a)^{*}T(b)p$ for
all $a,b,c$ in $A$;
2. $(b)$
$T(\cdot)p:A\to B^{**}$ is a triple homomorphism.
In the real setting, the lacking of a standard polarisation identity forces us
to make a slight modification of the argument given by Chu and Wong in [11,
Proposition 2.2], an sketch of the proof is included here for completeness.
###### Proposition 3.1.
Let $T:A\to B$ be a linear contraction between real C∗\- algebras. Then there
is a largest projection $p$ in $B^{\prime\prime}$ such that
1. $(a)$
$T\\{a,b,c\\}p=\\{T(a),T(b),T(c)\\}p$ and $pT(a)^{*}T(b)=T(a)^{*}T(b)p$ for
all $a,b,c$ in $A$;
2. $(b)$
$T(\cdot)p:A\to B^{\prime\prime}$ is a triple homomorphism.
Furthermore, the projection $p$ is closed.
###### Proof.
Let us define
$F_{1}:=\bigcap_{a,b\in
A_{1}}\Big{\\{}\varphi\in\mathcal{Q}(B):\|T(Q(a)b)-Q(T(a))(T(b))\|_{\varphi}=0\Big{\\}}$
$=\bigcap_{a,b\in
A_{1}}\Big{\\{}\varphi\in\mathcal{Q}(B):{\varphi}\Big{(}(\Theta_{a,b})(\Theta_{a,b})^{*}+(\Theta_{a,b})^{*}(\Theta_{a,b})\Big{)}=0\Big{\\}},$
where, for each $a,b\in A,$ $\Theta_{a,b}=T(Q(a)b)-Q(T(a))(T(b)).$ Clearly,
$F_{1}$ is a weak∗ closed face of $\mathcal{Q}(B)$ containing zero. Given
$a,b\in A_{1},$ we define a continuous affine mapping
$\Phi_{a,b}:\mathcal{Q}(B)\to\mathcal{Q}(B),$
$\varphi\mapsto\Phi_{a,b}(\varphi)(\cdot):=\varphi(T(b)^{*}T(a)\cdot
T(a)^{*}T(b)).$ For each natural $n\geq 2$,
$F_{n}:=\bigcap_{a,b\in
A_{1}}\Big{(}F_{n-1}\bigcap\Phi_{a,b}^{-1}(F_{n-1})\Big{)}$
is a weak∗ closed face of $\mathcal{Q}(B)$ containing zero and $F_{n}\subseteq
F_{n-1}.$ Therefore $F=\bigcap_{n=1}^{\infty}F_{n}$ is a weak∗ closed face of
$\mathcal{Q}(B)$ containing zero, and hence there exists a closed projection
$p\in B^{\prime\prime}$ satisfying that
$F=\\{\varphi\in\mathcal{Q}(B):\varphi(1-p)=0\\}$ (compare Lemma 2.6 or [32,
Theorem 5.2]). In particular, for each $a,b\in A$ and $\varphi\in F$ we have
$0=\Phi_{a,b}(\varphi)(1-p)=\varphi(T(b)^{*}T(a)(1-p)T(a)^{*}T(b))$
$=\varphi\Big{(}((1-p)T(a)^{*}T(b))^{*}((1-p)T(a)^{*}T(b))\Big{)}.$
Since the support tripotent of $\varphi$, $s(\varphi)$, is a projection with
$s(\varphi)\leq p\leq 1$, we have:
$\left\|(1-p)T(a)^{*}T(b)\right\|_{\varphi}^{2}=\varphi\left\\{(1-p)T(a)^{*}T(b),(1-p)T(a)^{*}T(b),p\right\\}$
$=\varphi P_{2}(p)\left\\{(1-p)T(a)^{*}T(b),(1-p)T(a)^{*}T(b),p\right\\}$
$=\varphi P^{1}(p)\left\\{(1-p)T(a)^{*}T(b),(1-p)T(a)^{*}T(b),p\right\\}$
$=\frac{1}{2}\varphi\Big{(}((1-p)T(a)^{*}T(b))((1-p)T(a)^{*}T(b))^{*}p\Big{)}$
$+\frac{1}{2}\varphi\Big{(}p((1-p)T(a)^{*}T(b))^{*}((1-p)T(a)^{*}T(b))\Big{)}$
$=\frac{1}{2}\varphi\Big{(}((1-p)T(a)^{*}T(b))^{*}((1-p)T(a)^{*}T(b))\Big{)}=0$
for every $a,b\in A$ and $\varphi\in F$. Observing that
$P_{2}(p)\left\\{(1-p)T(a)^{*}T(b),(1-p)T(a)^{*}T(b),p\right\\}$
$=P^{1}(p)\left\\{(1-p)T(a)^{*}T(b),(1-p)T(a)^{*}T(b),p\right\\}$
is a positive element in $(B^{\prime\prime})^{1}(p)$ (see $(\ref{equation
positive element})$) and $F=\\{\varphi\in
B^{\prime}:\varphi(p)=\varphi(1)=\|\varphi\|\leq 1\\}$ is the normal quasi-
state of the JBW-algebra $(B^{\prime\prime})^{1}(p)$, we have
$P_{2}(p)\left\\{(1-p)T(a)^{*}T(b),(1-p)T(a)^{*}T(b),p\right\\}=0$ for every
$a,b\in A$. This implies that
$p\Big{(}(1-p)T(a)^{*}T(b)T(b)^{*}T(a)(1-p)+pT(b)^{*}T(a)(1-p)T(a)^{*}T(b)\Big{)}^{*}p=0,$
or equivalently,
$\Big{(}(1-p)T(a)^{*}T(b)p\Big{)}^{*}\Big{(}(1-p)T(a)^{*}T(b)p\Big{)}=pT(b)^{*}T(a)(1-p)T(a)^{*}T(b)p=0,$
and hence
(3.1) ${(1-p)T(a)^{*}T(b)}p=0$
for every $a,b\in A$, which proves the second statement in $(a)$. By
definition of $F$, for each $a,b\in A$ and $\varphi\in F$, we have
${\varphi}\Big{(}(\Theta_{a,b})(\Theta_{a,b})^{*}+(\Theta_{a,b})^{*}(\Theta_{a,b})\Big{)}=0,$
which gives
$0={\varphi}P_{2}(p)\Big{(}(\Theta_{a,b})^{*}(\Theta_{a,b})\Big{)}={\varphi}P^{1}(p)\Big{(}(\Theta_{a,b})^{*}(\Theta_{a,b})\Big{)}.$
Recalling that $F$ coincides with the normal quasi-state of the JBW-algebra
$(B^{\prime\prime})^{1}(p)$ and
$P^{1}(p)\Big{(}(\Theta_{a,b})^{*}(\Theta_{a,b})\Big{)}=p(\Theta_{a,b})^{*}(\Theta_{a,b})p=(\Theta_{a,b}\
p)^{*}(\Theta_{a,b}\ p)$ is a positive element in $(B^{\prime\prime})^{1}(p)$,
we deduce that $\Theta_{a,b}\ p=0$, and accordingly
$T(\left\\{a,b,a\right\\})p=\left\\{T(a),T(b),T(a)\right\\}p.$
On the other hand,
$\left\\{T(a)p,T(b)p,T(a)p\right\\}=T(a)pT(b)^{*}T(a)p$ $=(\hbox{by }(\ref{eq
1 in prop linear
contraction}))=T(a)T(b)^{*}T(a)p=\left\\{T(a),T(b),T(a)\right\\}p,$
and we establish $(b)$.
Finally, when $q$ is a another projection in $B^{\prime\prime}$ satisfying
$(a)$ and $(b)$, it follows straightforwardly that
$F(q):=\\{\varphi\in\mathcal{Q}(B):\varphi(q)=\|\varphi\|\leq 1\\}$ is a norm-
closed face contained in $F_{1}$ and hence $F(q)\subseteq F$. It is known
that, under these conditions, $q\leq p$ (see, for example [32, Theorem 2.2]).
∎
The projection $p\in B^{\prime\prime}$ given by Proposition 3.1 is called the
_structure projection_ of $T$ and is denoted by $p_{T}$. Unfortunately, in
some cases, the structure projection $p_{T}$ satisfies that $T(.)p_{T}$
reduces to zero, even under the hypothesis of $T$ being an isometry. In the
complex setting, Chu and Wong proved in [11, Therorem 3.10] that $T(.)p_{T}$
is an isometry whenever $T:A\to B$ is an (complex) linear isometry between
C∗-algebras with $A$ abelian. We shall see next that the same statement is not
valid in the setting of real C∗-algebras.
###### Example 3.2.
Let $K$ be a compact Hausdorff space and let $C(K)_{\mathbb{R}}$ be the real
abelian C∗-algebra of complex continuous functions on $K$ regarded as a real
Banach space. Let $C(K,M_{2}(\mathbb{R}))$ be the real C∗-algebra of
continuous $M_{2}(\mathbb{R})$-valued functions on $K$, where
$M_{2}(\mathbb{R})$ denotes the real C∗-algebra of $2\times 2$ real matrices.
Define $T:C(K)_{\mathbb{R}}\rightarrow C(K,M_{2}(\mathbb{R}))$ by
$T(f)=\left(\begin{array}[]{ccc}\textrm{Re}(f)&\textrm{Im}(f)\\\ 0&0\\\
\end{array}\right)\;\;\;\;(f\in C(K)_{\mathbb{R}}).$
Then $T$ is a real linear isometry. We have $T(f)^{3}=T(f^{3})$ for all $f\in
C(K)_{\mathbb{R}}$, but in general $T(\\{f,g,h\\})\neq\\{T(f),T(g),T(h)\\}$.
We shall now show that the structure projection of $T$ satisfies
$p_{T}\leq\left(\begin{array}[]{cc}0&0\\\ 0&1\\\ \end{array}\right)$, and in
this case, $T(.)p_{T}=0$ (cf. [1]). Let $p$ be a projection satisfying the
thesis of Proposition 3.1. Then
$\displaystyle\\{T(f),T(g),T(f)\\}$
$\displaystyle=\left(\begin{array}[]{ccc}\textrm{Re}(f)&\textrm{Im}(f)\\\
0&0\\\ \end{array}\right)\left(\begin{array}[]{ccc}\textrm{Re}(g)&0\\\
\textrm{Im}(g)&0\\\
\end{array}\right)\left(\begin{array}[]{ccc}\textrm{Re}(f)&\textrm{Im}(f)\\\
0&0\\\ \end{array}\right)$
$\displaystyle=\left(\begin{array}[]{ccc}\textrm{Re}^{2}(f)\textrm{Re}(g)&\textrm{Im}^{2}(f)\textrm{Im}(g)\\\
+\textrm{Re}(f)\textrm{Im}(f)\textrm{Im}(g)&+\textrm{Re}(f)\textrm{Im}(f)\textrm{Re}(g)\\\
0&0\\\ \end{array}\right)$
and
$T(\\{f,g,f\\})=\left(\begin{array}[]{ccc}\textrm{Re}^{2}(f)\textrm{Re}(g)&\textrm{Im}^{2}(f)\textrm{Im}(g)\\\
-\textrm{Im}^{2}(f)\textrm{Re}(g)&-\textrm{Re}^{2}(f)\textrm{Im}(g)\\\
+2\textrm{Re}(f)\textrm{Im}(f)\textrm{Im}(g)&+2\textrm{Re}(f)\textrm{Im}(f)\textrm{Re}(g)\\\
0&0\\\ \end{array}\right).$
By assumptions:
$\displaystyle\left(\begin{array}[]{ccc}\textrm{Re}^{2}(f)\textrm{Re}(g)&\textrm{Im}^{2}(f)\textrm{Im}(g)\\\
+\textrm{Re}(f)\textrm{Im}(f)\textrm{Im}(g)&+\textrm{Re}(f)\textrm{Im}(f)\textrm{Re}(g)\\\
0&0\\\ \end{array}\right)p$
$\displaystyle=\left(\begin{array}[]{ccc}\textrm{Re}^{2}(f)\textrm{Re}(g)&\textrm{Im}^{2}(f)\textrm{Im}(g)\\\
-\textrm{Im}^{2}(f)\textrm{Re}(g)&-\textrm{Re}^{2}(f)\textrm{Im}(g)\\\
+2\textrm{Re}(f)\textrm{Im}(f)\textrm{Im}(g)&+2\textrm{Re}(f)\textrm{Im}(f)\textrm{Re}(g)\\\
0&0\\\ \end{array}\right)p$
which simplifies to
(3.2)
$\left(\begin{array}[]{ccc}\textrm{Re}(f)\textrm{Im}(f)\textrm{Im}(g)&\textrm{Re}(f)\textrm{Im}(f)\textrm{Re}(g)\\\
-\textrm{Im}^{2}(f)\textrm{Re}(g)&-\textrm{Re}^{2}(f)\textrm{Im}(g)\\\ 0&0\\\
\end{array}\right)p=0.$
Let $f,g\in C(K)_{\mathbb{R}}$ be the constant functions $g(x)=1$ and
$f(x)=\imath$. Substituting $f$ and $g$ into (3.2) we obtain
$\left(\begin{array}[]{ccc}-1&0\\\ 0&0\\\ \end{array}\right)p(x)=0\qquad(x\in
K),$
which gives the desired statement.
According to the result established by Chu and Wong in [11, Therorem 3.10],
the mapping $T$ cannot be complexified to an isometry on the complexification
$C(K)_{\mathbb{R}}\oplus\imath C(K)_{\mathbb{R}}$. Here is a direct argument:
let $T_{c}$ denote the complexification of $T$.
$\displaystyle T_{c}(g+if)$ $\displaystyle=T(g)+iT(f)$
$\displaystyle=\left(\begin{array}[]{ccc}\textrm{Re}(g)+i\textrm{Re}(f)&\textrm{Im}(g)+i\textrm{Im}(f)\\\
0&0\\\ \end{array}\right)\in C(K,M_{2}(\mathbb{C})).$
Let $g(x)=1$ and $f(x)=\imath$ be the constant functions considered before.
Then $\|g+if\|=2$ while
$\|T_{c}(g+if)\|=\left\|\left(\begin{array}[]{ccc}1&i\\\ 0&0\\\
\end{array}\right)\right\|_{M_{2}(\mathbb{R})_{c}}=\sqrt{2}.$
Therefore $T_{c}:C(K)_{\mathbb{R}}\bigoplus\imath C(K)_{\mathbb{R}}\rightarrow
C(K,M_{2}(\mathbb{C}))$ is not an isometry.
We shall see later, that the obstacle in the above example relies on the
“complex nature” of $C(K)_{\mathbb{R}}$ and the real linearity of $T$.
Let $A$ be a real C∗-algebra, a _“complex” character_ of $A$ is a real linear
homomorphism $\rho:A\to\mathbb{C}$. Since every complex character
$\rho:A\to\mathbb{C}$ admits a complex linear extension to a character
$\widetilde{\rho}:A_{c}\to\mathbb{C}$, we may assume, via Gelfand theory, that
$\rho$ is a real linear ∗-homomorphism. Every non-zero complex character of
$A$ is unital whenever $A$ has a unit element, or can be extended to a unital
complex character of the unitisation of $A$ otherwise. Following standard
notation, $\Omega(A)$ will denote the space of all complex characters of $A$.
We write
$\Omega(A)_{\mathbb{F}}:=\\{\rho\in\Omega(A):\rho(A)=\mathbb{F}\\},\hbox{
where }\mathbb{F}=\mathbb{R}\hbox{ or }\mathbb{C}.$
The elements in $\Omega(A)_{\mathbb{R}}$ (resp., in $\Omega(A)_{\mathbb{C}}$)
are called _complex characters of real type_ (resp., _complex characters of
complex type_). When $A$ is abelian, we shall say that $A$ is of _real type_
if $\Omega(A)_{\mathbb{C}}=\emptyset$, and of _complex type_ if
$\Omega(A)_{\mathbb{C}}=\emptyset$. An abelian real C∗-algebra $A$ is of real
type (resp., of complex type) if and only if all non-zero complex characters
of $A$ are onto $\mathbb{R}$ (resp., onto $\mathbb{C}$) (compare [29, Theorem
2.7.7]). In the more general setting of real commutative J∗-algebras, real and
complex types where studied by L.J. Bunce and C.-H. Chu in [4, §3]. When $A$
is an abelian complex C∗-algebra, then the real C∗-algebra underlying $A$,
$A_{\mathbb{R}}$, is a real C∗-algebra of complex type. Given a locally
compact Hausdorff space $X$, the real C∗-algebra $C_{0}(X,\mathbb{R})$ is of
real type. Furthermore, it follows from [4, Lemma 3.4] that every abelian real
C∗-algebra of real type if of the form $C_{0}(X,\mathbb{R})$, for some locally
compact Hausdorff space $X$.
Our next result shows that the main result in [11] can be also proved when the
domain is an abelian real C∗-algebra of real type.
Henceforth, the closed unit ball of a Banach space $E$ will be denoted by
$E_{1}$ and the set of extreme points of a convex set $S$ will be denoted by
$\partial S$.
###### Proposition 3.3.
Let $T:X\to Y$ be a (not necessarily surjective) linear isometry between
Banach spaces and let $E=T(X)$. For each $\psi\in\partial E^{\prime}_{1},$ the
set
$\mathcal{Q}_{\psi}=\left\\{\phi\in\partial
Y^{\prime}_{1}:\phi|_{E}=\psi\right\\}$
is non-empty, that is, there exists $\phi\in\partial Y^{\prime}_{1}$
satisfying $\psi=\phi|_{E}$ and $T^{\prime}(\phi)\in\partial X^{\prime}_{1}$.
###### Proof.
Let $T^{\prime}:E^{\prime}\rightarrow X^{\prime}$ be the dual map of the
surjective isometry $T:X\rightarrow E$. Then $T^{\prime}$ is a linear
surjective isometry. Let $\psi$ be an extreme point of the unit ball,
$E_{1}^{\prime}$, of $E^{\prime}$. By the Hahn-Banach theorem, the set
$\mathcal{F}_{\psi}=\left\\{\phi\in Y^{\prime}_{1}:\phi|_{E}=\psi\right\\}$
is a non-empty weak∗ closed face of $Y^{\prime}_{1}$, thus, by the Krein-
Milman theorem, there exists $\phi\in\mathcal{F}_{\psi}\bigcap\partial
Y^{\prime}_{1}=\mathcal{Q}_{\psi}$. By an abuse of notation, we also denote by
$T^{\prime}:E^{\prime}\to X^{\prime}$ the dual map of $T:X\to Y$ since
confusion is unlikely. One sees readily that
$T^{\prime}(\phi)=T^{\prime}(\psi)$. Since $\psi=\phi|_{E}$ is an extreme
point in $E_{1}^{\prime}$, $T^{\prime}(\psi)$ is an extreme point in
$X^{\prime}_{1}$ because $T^{\prime}:E^{\prime}\rightarrow X^{\prime}$ is a
surjective linear isometry. ∎
We show below that if a real C∗-algebra $A$ admits a complex character of real
type (i.e. $\Omega(A)_{\mathbb{R}}\neq\emptyset$) then the mapping
$T(\cdot)p_{T}$ is non-zero for every linear isometry from $A$ into another
C∗-algebra. Our next result is in fact an appropriate real version of [10,
Proposition 2] and [11, Proposition 4.3].
###### Theorem 3.4.
Let $T:A\to B$ be a (not necessarily surjective) linear isometry between real
C∗-algebras. Suppose $\rho\in\Omega_{\mathbb{R}}(A)$, then there exists
$\phi\in\partial B^{\prime}_{1}$ and a minimal partial isometry $u_{\phi}\in
B^{\prime\prime}$ such that $T^{\prime}(\phi)=\rho$, $u_{\phi}$ is the support
partial isometry of $\phi$,
${T^{\prime\prime}(x)}=\rho(x){u_{\phi}}+P_{0}(u_{\phi})(T^{\prime\prime}(x)),$
$\\{u_{\phi},u_{\phi},\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\}\\}=\\{u_{\phi},u_{\phi},T^{\prime\prime}(\\{x,y,z\\})\\},$
and
$\\{u_{\phi},\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\},u_{\phi}\\}=\\{u_{\phi},T^{\prime\prime}(\\{x,y,z\\}),u_{\phi}\\},$
for every $x,y,z\in A^{\prime\prime}.$ Moreover, the minimal projection
$p_{\phi}=u_{\phi}^{*}u_{\phi}$ satisfies statements $(a)$ and $(b)$ in
Proposition 3.1 and the mapping $T(\cdot)p_{\phi}:A\to B^{\prime\prime}$ is a
non-zero triple homomorphism. In particular, $T(\cdot)p_{T}$ is non-zero.
###### Proof.
Let $\rho\in\Omega(A)_{\mathbb{R}}$, where $A$ is a (not necessarily abelian)
real C∗-algebra. The map $\rho^{\prime\prime}:A^{\prime\prime}\to\mathbb{R}$
is a surjective norm-one weak∗ continuous *-homomorphism. The set
$K=\ker(\rho^{\prime\prime})$ is a weak∗ closed ideal of $A^{\prime\prime}$,
thus there exists a minimal and reduced central projection $p\in
A^{\prime\prime}$ such that $A^{\prime\prime}=K\oplus^{\infty}\mathbb{R}p$ and
$K=A^{\prime\prime}(1-p).$ In particular, $\rho\in\partial A^{\prime}_{1}.$
If $T:A\to B$ is a linear isometry, taking $E=T(A),$ $T^{\prime}:E^{\prime}\to
A^{\prime}$ is a surjective linear isometry, then there exists
$\psi\in\partial E^{\prime}_{1}$ such that $T^{\prime}(\psi)=\rho$.
Proposition 3.3 assures the existence of $\phi\in\partial B^{\prime}_{1}$ such
that $T^{\prime}(\phi)=T^{\prime}(\psi)=\rho$. We observe that
$(A^{\prime\prime})^{-1}(p)=\\{0\\},$ $P^{1}(p)(x)=xp$ and
$T^{\prime\prime}:A^{\prime\prime}\to B^{\prime\prime}$ is a weak* continuous
linear isometry.
By the remarks following Lemma 2.3 (cf. [37, Corollary 2.1]), the minimal
tripotents of $B^{\prime\prime}$ are the support tripotents of the extreme
points of $B^{\prime}_{1}$. Let $u_{\phi}\in B^{\prime\prime}$ be the support
tripotent of $\phi$. Then $\phi(.)=\phi Q(u_{\phi})(.)$ and
$P^{1}(u_{\phi})(.)=\phi(.)u_{\phi}$ (cf. [37, Lemma 2.7]). Thus, $\phi
T(x)=T^{\prime}(\phi)(x)=\rho(x)$ implies
(3.3) $P^{1}(u_{\phi})T(x)=\phi(T(x))u_{\phi}=\rho(x)u_{\phi},$
for every $x\in A$. Since $T^{\prime\prime}:A^{\prime\prime}\to
B^{\prime\prime}$ and $\rho\in A^{\prime}$ are weak* continuous and $A_{1}$ is
weak* dense in $A^{\prime\prime}_{1}$, we have
(3.4)
$P^{1}(u_{\phi})T^{\prime\prime}(x)=\phi(T^{\prime\prime}(x))u_{\phi}=\rho(x)u_{\phi},$
for every $x\in A^{\prime\prime}$. In particular,
$P^{1}(u_{\phi})T^{\prime\prime}(\\{x,y,z\\})=\rho\left\\{x,y,z\right\\}u_{\phi}=\rho(x)\rho(y)\rho(z)\\{u_{\phi},u_{\phi},u_{\phi}\\}$
$=\\{P^{1}(u_{\phi})T^{\prime\prime}(x),P^{1}(u_{\phi})T^{\prime\prime}(y),P^{1}(u_{\phi})T^{\prime\prime}(z)\\},$
for every $x,y,z\in A^{\prime\prime}$, that is,
$P^{1}(u_{\phi})T^{\prime\prime}:A^{\prime\prime}\to B^{\prime\prime}$ is a
non-zero triple homomorphism.
We shall show that $\\{u_{\phi},u_{\phi},T(a)\\}=u_{\phi}$ for every $a\in
A^{\prime\prime}$ satisfying $\|a\|=1$ and $P^{1}(p)(a)=p$. By $(\ref{correct
eq 3.1b})$ and the Cauchy-Schwarz inequality (cf. Lemma 2.4), we have
$\displaystyle 1$ $\displaystyle=\|P^{1}(p)(a)\|=|\phi\circ
T^{\prime\prime}(a)|^{2}=|\phi(\\{u_{\phi},T^{\prime\prime}(a),u_{\phi}\\})|^{2}$
$\displaystyle\leq\phi(\\{u_{\phi},u_{\phi},u_{\phi}\\})\phi(\\{T^{\prime\prime}(a),T^{\prime\prime}(a),u_{\phi}\\})$
$\displaystyle\leq\|T^{\prime\prime}(a)\|^{2}=\|a\|^{2}=1$
which implies
$\phi(\\{T^{\prime\prime}(a),T^{\prime\prime}(a),u_{\phi}\\})=1.$
Let $N_{\phi}=\\{b\in B^{\prime\prime}\;:\;\phi(\\{b,b,u_{\phi}\\})=0\\}$. By
Lemma 2.5,
(3.5) $N_{\phi}=P_{0}(u_{\phi})(B^{\prime\prime}).$
Now we claim that $T^{\prime\prime}(a)-u_{\phi}\in N_{\phi}$, for any $a$ as
above. Indeed, taking into account that, by Lemma 2.4,
$\phi(\\{T^{\prime\prime}(a),u_{\phi},u_{\phi}\\})=\phi(\\{u_{\phi},T^{\prime\prime}(a),u_{\phi}\\})$,
we have
$\displaystyle\phi(\\{T^{\prime\prime}(a)-u_{\phi},T^{\prime\prime}(a)-u_{\phi},u_{\phi}\\})$
$\displaystyle=\phi(\\{T^{\prime\prime}(a),T^{\prime\prime}(a),u_{\phi}\\})-2\phi(\\{u_{\phi},T^{\prime\prime}(a),u_{\phi}\\})+\phi(u_{\phi})$
$\displaystyle=\phi(\\{T^{\prime\prime}(a),T^{\prime\prime}(a),u_{\phi}\\})-2\phi(T^{\prime\prime}(a))+\phi(u_{\phi})$
$\displaystyle=1-2\rho(a)+1=0.$
Therefore, by (3.5), $T^{\prime\prime}(a)-u_{\phi}\in
P_{0}(u_{\phi})(B^{\prime\prime})$ and then
(3.6) $\\{u_{\phi},u_{\phi},T^{\prime\prime}(a)\\}=u_{\phi},$
for every $a$ as above.
In the next step we prove that
$\phi(\\{T^{\prime\prime}(b),T^{\prime\prime}(b),u_{\phi}\\})=0$ whenever
$b\in A^{\prime\prime}$ satisfies $P^{1}(p)(b)=0$. Without loss of generality
we can take $\|b\|=1$. Since
$A^{\prime\prime}=A^{\prime\prime}(1-p)\oplus^{\infty}\mathbb{R}p$, $p$ and
$b$ are orthogonal in $A^{\prime\prime}$, $\|p+b\|=1$ and $P^{1}(p)(b+p)=p$.
We deduce, from the above arguments, that
$T^{\prime\prime}(b+p)+N_{\phi}=u_{\phi}+N_{\phi}=T^{\prime\prime}(p)+N_{\phi},$
which implies $T^{\prime\prime}(b)\in N_{\phi}$ and
$\phi(\\{T^{\prime\prime}(b),T^{\prime\prime}(b),u_{\phi}\\})=0$.
Now, let $c\in A^{\prime\prime}$ with $\|c\|=1$. Then $P^{1}(p)(c-\rho(c)p)=0$
and, by the arguments of the previous paragraph, we have
$T^{\prime\prime}(c-\rho(c)p)\in N_{\phi}$. That is,
$\\{u_{\phi},u_{\phi},T^{\prime\prime}(c-\rho(c)p)\\}=0$ and hence, by
$(\ref{eq new theorem 1})$,
(3.7)
$\\{u_{\phi},u_{\phi},T^{\prime\prime}(c)\\}=\rho(c)\\{u_{\phi},u_{\phi},T^{\prime\prime}(p)\\}=\rho(c)u_{\phi}.$
Left multiply by $u_{\phi}u_{\phi}^{*}$, right multiply by
$u_{\phi}^{*}u_{\phi}$ and subtracting, we have
$u_{\phi}u_{\phi}^{*}T^{\prime\prime}(c)=T^{\prime\prime}(c)u_{\phi}^{*}u_{\phi}=u_{\phi}u_{\phi}^{*}T^{\prime\prime}(c)u_{\phi}^{*}u_{\phi}$,
for every $c\in A^{\prime\prime}$.
Therefore
(3.8)
${T^{\prime\prime}(c)}=\rho(c){u_{\phi}}+P_{0}(u_{\phi})(T^{\prime\prime}(c)),$
for every $c\in A^{\prime\prime}$. It follows that
$T^{\prime\prime}(\\{x,y,z\\})=\rho\\{x,y,z\\}\
u_{\phi}+P_{0}(u_{\phi})(T^{\prime\prime}(\\{x,y,z\\})),$
and by Peirce rules
$\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\}=$
$=\rho(x)\rho(y)\rho(z)u_{\phi}+\\{{P_{0}(u_{\phi})(T^{\prime\prime}(x))},{P_{0}(u_{\phi})(T^{\prime\prime}(y))},{P_{0}(u_{\phi})(T^{\prime\prime}(z))}\\}$
$=\rho\\{x,y,z\\}u_{\phi}+\\{{P_{0}(u_{\phi})(T^{\prime\prime}(x))},{P_{0}(u_{\phi})(T^{\prime\prime}(y))},{P_{0}(u_{\phi})(T^{\prime\prime}(z))}\\},$
which assures that
(3.9)
$\\{u_{\phi},u_{\phi},\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\}\\}=\\{u_{\phi},u_{\phi},T^{\prime\prime}(\\{x,y,z\\})\\},$
and
(3.10)
$\\{u_{\phi},\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\},u_{\phi}\\}=\\{u_{\phi},T^{\prime\prime}(\\{x,y,z\\}),u_{\phi}\\},$
for every $x,y,z\in A^{\prime\prime}.$
Let $p_{\phi}=u_{\phi}^{*}u_{\phi}$ and $q_{\phi}=u_{\phi}u_{\phi}^{*}$ be the
initial and final minimal projections of the minimal tripotent $u_{\phi}$ (see
Lemma 2.1). We then have
$(\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\}-T^{\prime\prime}(\\{x,y,z\\}))p_{\phi}+$
$q_{\phi}(\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\}-T^{\prime\prime}(\\{x,y,z\\}))=0,$
and
$u_{\phi}(\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\}-T^{\prime\prime}(\\{x,y,z\\}))^{*}u_{\phi}=0,$
and consequently:
$q_{\phi}(\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\}-T^{\prime\prime}(\\{x,y,z\\}))p_{\phi}=0,$
and
(3.11)
$(\\{T^{\prime\prime}(x),T^{\prime\prime}(y),T^{\prime\prime}(z)\\}-T^{\prime\prime}(\\{x,y,z\\}))p_{\phi}=0,$
witnessing that $T^{\prime\prime}(\cdot)p_{\phi}:A\to B^{\prime\prime}$ is a
non-zero triple homomorphism (observe that
$T^{\prime\prime}(\cdot)p_{\phi}=\rho(\cdot)u_{\phi}\neq 0$). Proposition 3.1
implies that $p\leq p_{T}$ and $T(\cdot)p_{T}\neq 0$. ∎
Suppose $A$ is an abelian real C∗-algebra of complex type, that is,
$\Omega(A)_{\mathbb{R}}=\emptyset$. In general, $A$ need not be C∗-isomorphic
to $C_{0}(X)_{\mathbb{R}}$ (compare, for example, [4, Remark 3.8]). However,
by [4, Theorem 3.7], $A^{\prime\prime}$ is isometric and C∗-isomorphic to
$C(\Omega)_{\mathbb{R}}$, for some compact hyperstonean space $\Omega$. Let
$J_{A}:A\hookrightarrow A^{\prime\prime}$ denote the canonical inclusion of
$A$ into $A^{\prime\prime}$, then $J_{A}$ is an isometric C∗-embedding. Let
$T:A^{\prime\prime}=C(\Omega)_{\mathbb{R}}\rightarrow
C(\Omega,M_{2}(\mathbb{R}))$ be the isometry given in Example 3.2. That is
$T(f)=\left(\begin{array}[]{ccc}\textrm{Re}(f)&\textrm{Im}(f)\\\ 0&0\\\
\end{array}\right)\;\;\;\;(f\in C(\Omega)_{\mathbb{R}}).$
The mapping $S=TJ_{A}:A\to C(\Omega,M_{2}(\mathbb{R}))$ is an isometry
satisfying $S(.)p_{S}=0$.
When $T$ is an isometry from an abelian real C∗-algebra of real type to
another real C∗-algebra, we can prove that $T(.)p_{T}$ is an isometry.
###### Theorem 3.5.
Let $B$ be a real C∗-algebra and $T:C_{0}(X,\mathbb{R})\rightarrow B$ be a
(not necessarily surjective) linear isometry. Let $p_{T}\in B^{\prime\prime}$
be the structure projection of $T$. Then
$T(.)p_{T}:C_{0}(X,\mathbb{R})\rightarrow B^{\prime\prime}$ is an isometric
triple homomorphism.
###### Proof.
Let us denote $A=C_{0}(X,\mathbb{R})$ and $E=T(A)$ and let
$T^{\prime}:E^{\prime}\rightarrow A^{\prime}$ be the dual map of the
surjective isometry $T:A\rightarrow E$. The real C∗-algebra $A$ is of real
type and $\Omega(A)_{\mathbb{R}}=\\{\delta_{x}:x\in X\\}$. By Proposition 3.3,
the set
$\mathcal{Q}=\left\\{\phi\in\partial
B^{\prime}_{1}:\phi|_{E}\in\Omega(A)_{\mathbb{R}}\right\\}$
is non-empty. By Theorem 3.4, for each $x\in X$, there exists
$\phi_{x}\in\mathcal{Q}$ and a minimal partial isometry $u_{x}\in
B^{\prime\prime}$ such that $T^{\prime}(\phi_{x})=\delta_{x}$, $u_{x}$ is the
support partial isometry of $\phi_{x}$,
(3.12) ${T(f)}=\delta_{x}(f){u_{x}}+P_{0}(u_{x})(T(f)),$
$(\\{T(f),T(g),T(h)\\}-T(\\{f,g,h\\}))p_{x}=0,$
and
$p_{x}T(f)^{*}T(g)=T(f)^{*}T(g)p_{x},$
for every $f,g,h\in A$, where $p_{x}=u_{x}^{*}u_{x}$ is a minimal projection
in $B^{\prime\prime}$. Let $p=\bigvee_{x\in X}p_{x}$ be the lattice supremum
in $B^{\prime\prime}$. By Lemma 2.3 and (3.12), we obtain
(3.13) $\\{T(f),T(g),T(h)\\}p=T(\\{f,g,h\\})p\hbox{ and
}pT(f)^{*}T(g)=T(f)^{*}T(g)p,$
for every $f,g,h\in A$. Thus, $T(\cdot)p:A\to B^{\prime\prime}$ is a triple
homomorphism. Proposition 3.1 implies $p\leq p_{T}.$
We claim that $T(\cdot)p$ is an isometry. Indeed, for each $f\in A$, we can
pick $x_{0}\in X$ with $\|f\|=|f(x_{0})|$. Let $\phi_{x_{0}}\in\mathcal{Q}$ be
such that $T^{\prime}(\phi_{x_{0}})=\delta_{x_{0}}$. In this case,
$\|f\|=\|T(f)\|\geq\|T(f)p\|\geq\|T(f)pp_{x_{0}}\|=\|T(f)p_{x_{0}}\|$
$=\|T(f)u_{x_{0}}^{*}u_{x_{0}}\|\geq\|u_{x_{0}}u_{x_{0}}^{*}T(f)u_{x_{0}}^{*}u_{x_{0}}\|$
$=\|\\{u_{x_{0}},\\{u_{x_{0}},T(f),u_{x_{0}}\\},u_{x_{0}}\\}\|=\|f(x_{x_{0}})u_{x_{0}}\|=\|f\|,$
which implies $\|T(f)p\|=\|f\|$ for all $f\in A$.
Finally, since $p\leq p_{T}$ we have
$\|f\|=\|T(f)p\|=\|T(f)p_{T}p\|\leq\|T(f)p_{T}\|\leq\|T(f)\|=\|f\|,$
which completes the proof. ∎
Let $E$ be a real or complex JBW∗-triple, two elements $\phi,\psi\in
E_{{}_{{}^{\prime}}}$ are said to be _orthogonal_ (written $\phi\perp\psi$) if
their support tripotents are orthogonal, that is, $u_{\phi}\perp u_{\psi}$.
C.M. Edwards and G. Rüttimann proved in [16, Theorem 5.4] that in the setting
of complex JBW∗-triples, $\phi$ and $\psi$ are orthogonal if, and only if,
$\|\phi\pm\psi\|=\|\phi\|+\|\psi\|$.
Suppose $E$ is a real JBW∗-triple whose complexification, $E_{c}$, is a
(complex) JBW∗-triple. We have already commented that there is a conjugate
linear isometry $\tau:E_{c}\rightarrow E_{c}$ of period 2 such that
$E=E_{c}^{\tau}=\\{b\in E_{c}\>:\>\tau(b)=b\\}$. The “dual” map
$\widetilde{\tau}:E_{c}{{}^{*}}\rightarrow E_{c}^{*}$ defined by
$\widetilde{\tau}(\phi)(b)=\overline{\phi(\tau(b))}\qquad(\phi\in
E_{c}^{*},\,b\in E_{c})$
is a conjugate linear isometry of period 2 and the mapping
$(E_{c}^{*})^{\widetilde{\tau}}\to E^{\prime}$ $\varphi\mapsto\varphi|_{E}$
is a surjective linear isometry. Moreover, by [31, Proposition 2.3], $\tau$ is
weak* continuous and hence $\widetilde{\tau}((E_{c})_{*})=(E_{c})_{*}$ and the
restricted mapping
$((E_{c})_{*})^{\widetilde{\tau}}\to E_{{}_{{}^{\prime}}}$
$\varphi\mapsto\varphi|_{E}$
is a surjective linear isometry. Take two elements $\phi,\psi\in
E_{{}_{{}^{\prime}}}$, then $\phi\perp\psi$ (in $E_{{}_{{}^{\prime}}}$) if,
and only if, $\phi\perp\psi$ as elements in $E_{*}$, therefore
$\|\phi\pm\psi\|=\|\phi\|+\|\psi\|$ as elements in $E_{*}$, which is
equivalent to $\|\phi\pm\psi\|=\|\phi\|+\|\psi\|$ as elements in
$E_{{}_{{}^{\prime}}}$.
Let $\mathcal{U}(E)$ and $\mathcal{U}(E_{c})$ denote the sets of all
tripotents in $E$ and $E_{c}$, respectively. Let $\mathcal{U}(E_{c})^{\sim}$
denote the union of the set $\mathcal{U}(E_{c})$ and a one point set
$\\{\omega\\}$ such that $e\leq\omega$ for every $e\in\mathcal{U}(E_{c})$. It
is known that $\mathcal{U}(E_{c})^{\sim}$ is a complete lattice and
$\mathcal{U}(E)\cup\\{\omega\\}=\\{e\in\mathcal{U}(E_{c}):\tau(e)=e\\}\cup\\{\omega\\}$
is a sub-complete lattice of $\mathcal{U}(E_{c})^{\sim}$. The supremum of a
family $(e_{i})\subset\mathcal{U}(E_{c})$ (resp., in $\mathcal{U}(E)$) need
not exist, in general, in $\mathcal{U}(E_{c})$ (resp., in $\mathcal{U}(E)$).
However, for every family $(e_{i})$ of mutually orthogonal elements in
$\mathcal{U}(E_{c})$ (resp., in $\mathcal{U}(E)$), the supremum
$\bigvee_{i}e_{i}$ exists in $\mathcal{U}(E_{c})$ or in (resp., in
$\mathcal{U}(E)$) (cf. [16, Theorem 5.1] and [17]). Moreover, suppose that an
element $a\in E$ satisfies $a\perp e_{i}$ for every $i$, then
$a\perp\bigvee_{i}e_{i}$.
The next lemma subsumes some of the above results.
###### Lemma 3.6.
Let $E$ be a real JBW∗-triple. The following statements hold:
1. $(a)$
Let $\phi$ and $\psi$ be two elements in $E_{{}_{{}^{\prime}}}$, then
$\phi\perp\psi$ if, and only if, $\|\phi\pm\psi\|=\|\phi\|+\|\psi\|$.
2. $(b)$
Let $(e_{i})$ be a family of mutually orthogonal elements in $\mathcal{U}(E)$.
Then the supremum $\bigvee_{i}e_{i}$ exists in $\mathcal{U}(E)$. Moreover,
suppose $a\in E$ satisfies $a\perp e_{i}$, for every $i$, then
$\bigvee_{i}e_{i}\perp a$. $\hfill\Box$
The next result reveals the connection between the real versions of Theorem 1
in [10] and Theorem 3.10 in [11].
###### Theorem 3.7.
Let $B$ be a real JB∗-triple and $T:C_{0}(X,\mathbb{R})\rightarrow B$ be a not
necessarily surjective linear isometry. Then there exists a partial isometry
$u\in B^{\prime\prime}$ such that
${T(f)}=\delta_{x}(f){u_{x}}+P_{0}(u_{x})(T(f)),$
$\\{u,T(\\{f,g,h\\}),u\\}=\\{u,\\{T(f),T(g),T(h)\\},u\\},$
for every $f,g,h\in C_{0}(X,\mathbb{R})$, and
$\\{u,T(\cdot),u\\}:C(X,\mathbb{R})\rightarrow B^{\prime\prime}$ is an
isometry.
###### Proof.
Keeping in mind the notation in the proof of Theorem 3.5, we write
$A=C_{0}(X,\mathbb{R})$, $\Omega(A)_{\mathbb{R}}=\\{\delta_{x}:x\in X\\}$ and
$\mathcal{Q}=\left\\{\phi\in\partial
B^{\prime}_{1}:\phi|_{E}\in\Omega(A)_{\mathbb{R}}\right\\}.$
By Theorem 3.4, whose proof is valid when $B$ is a real JB∗-triple, for each
$x\in X$, there exists $\phi_{x}\in\mathcal{Q}$ and a minimal partial isometry
$u_{x}\in B^{\prime\prime}$ such that $T^{\prime}(\phi_{x})=\delta_{x}$,
$u_{x}$ is the support partial isometry of $\phi_{x}$,
(3.14) ${T(f)}=\delta_{x}(f){u_{x}}+P_{0}(u_{x})(T(f)),$
and hence
$T(\\{f,g,h\\})-\\{T(f),T(g),T(h)\\}\in B^{\prime\prime}_{0}(u_{x})$
for every $f,g,h\in A$ and $x\in X$.
We shall prove now that $u_{x}\perp u_{y}$ whenever $x\neq y$ in $X$. Indeed,
for $x\neq y$ we have $\delta_{x}\perp\delta_{y}$, and thus
$2=\|\delta_{x}\pm\delta_{y}\|=\|T^{\prime}(\phi_{x}\pm\phi_{y})\|\leq\|\phi_{x}\pm\phi_{y}\|\leq\|\phi_{x}\|+\|\phi_{y}\|=2.$
It follows from Lemma 3.6$(a)$ and the comments preceding it, that
$\phi_{x}\perp\phi_{y}$, or equivalently, $u_{x}\perp u_{y}$. Therefore
$(u_{x})_{x\in X}$ is a family of mutually orthogonal minimal tripotents in
$B^{\prime\prime}$. The supremum $u=\bigvee_{x}u_{x}$ exists and defines a
tripotent in $B^{\prime\prime}$ (cf. Lemma 3.6$(b)$). Since, by
$(\ref{4.3.13c})$, for each $f,g,$ and $h$ in $A$ and $x\in X$, we have
$T(\\{f,g,h\\})-\\{T(f),T(g),T(h)\\}\perp u_{x}$, we deduce from Lemma
3.6$(b)$, that $T(\\{f,g,h\\})-\\{T(f),T(g),T(h)\\}\perp u,$ for every
$f,g,h\in A$.
In order to prove the last statement, observe that for each $f\in A$ we can
take $x_{0}\in X$ with $\|f\|=|f(x_{0})|$. Let $\phi_{x_{0}}\in\partial
B^{\prime}_{1}$ such that $T^{\prime}(\phi_{x_{0}})=\delta_{x_{0}}$. Then,
$\displaystyle\|f\|=\|T(f)\|$
$\displaystyle\geq\|\\{u,T(f),u\\}\|\geq\|\\{u_{0},\\{u,T(f),u\\},u_{0}\\}\|$
$\displaystyle\geq\|Q(u_{x_{0}})Q(u)(T(f))\|=\|Q(u_{x_{0}})^{2}(T(f))\|$
$\displaystyle=\|\\{u_{0},\\{u_{0},T(f),u_{0}\\},u_{0}\\}\|=\|f(x_{x_{0}})u_{x_{0}}\|=\|f\|,$
which proves
$\|\\{u,T(f),u\\}\|=\|f\|\qquad(f\in A).$
∎
###### Remark 3.8.
When in the above Theorem 3.7, $B$ is a real C∗-algebra, the projection
$p=u^{*}u$ coincides with the one considered in the proof of Theorem 3.5.
Using the previous results, we can now show that if $T$ is an isometry between
two real C∗-algebras $A$ and $B$, then $T$ is a local triple homomorphism via
a tripotent in $B^{\prime\prime}$. That does not however imply that $T$ is a
triple homomorphism on the whole algebra $A$ (compare Example 3.2). The
following theorem is an extension of the results in [10, 11].
We recall that given an element $a$ in a real JB∗-triple $A$, the real
JB∗-subtriple, $A_{a},$ generated by $a$ is linearly isometric to the real
JB∗-triple $C_{0}(X,\mathbb{R})$ of real-valued continuous functions on $X$
vanishing at infinity (see [7] and [27]).
###### Theorem 3.9.
Let $T:A\rightarrow B$ be a not necessarily surjective linear isometry between
two real JB∗-triples. Then, for each $a\in A$, there exists a tripotent $u\in
B^{\prime\prime}$ such that
1. $(a)$
$\\{u,T(\\{f,g,h\\}),u\\}=\\{u,\\{T(f),T(g),T(h)\\},u\\}$, for all $f,g,h$ in
the real JB∗-subtriple generated by $a$;
2. $(b)$
The mapping $\\{u,T(\cdot),u\\}:A_{a}\rightarrow B^{\prime\prime}$ is a linear
isometry.
Furthermore, when $B$ is a real C∗-algebra, the projection $p=u^{*}u$
satisfies that $T(\cdot)p:A_{a}\rightarrow B^{\prime\prime}$ is an isometric
triple homomorphism.
###### Proof.
Since $A_{a}$ is linearly isometric to the JB∗-triple $C_{0}(X,\mathbb{R})$ of
all real-valued continuous functions on $X$ vanishing at infinity (cf. [27]),
the result follows immediately from Theorem 3.5, Theorem 3.7 and Remark 3.8. ∎
Acknowledgements: The useful comments and suggestions made by the Referee in
his/her report were fundamental to improve the presentation and quality of the
paper. We are grateful to the referee for his/her thorough suggestions.
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|
arxiv-papers
| 2013-09-16T07:27:52 |
2024-09-04T02:49:50.941845
|
{
"license": "Public Domain",
"authors": "Maria Apazoglou and Antonio M. Peralta",
"submitter": "Antonio M. Peralta",
"url": "https://arxiv.org/abs/1309.3838"
}
|
1309.3839
|
# Orthogonal forms and orthogonality preservers on real function algebras
Jorge J. Garcés [email protected] Departamento de Análisis Matemático, Facultad
de Ciencias, Universidad de Granada, 18071 Granada, Spain. and Antonio M.
Peralta [email protected] Departamento de Análisis Matemático, Facultad de
Ciencias, Universidad de Granada, 18071 Granada, Spain.
###### Abstract.
We initiate the study of orthogonal forms on a real C∗-algebra. Motivated by
previous contributions, due to Ylinen, Jajte, Paszkiewicz and Goldstein, we
prove that for every continuous orthogonal form $V$ on a commutative real
C∗-algebra, $A$, there exist functionals $\varphi_{1}$ and $\varphi_{2}$ in
$A^{*}$ satisfying
$V(x,y)=\varphi_{1}(xy)+\varphi_{2}(xy^{*}),$
for every $x,y$ in $A$. We describe the general form of a (not-necessarily
continuous) orthogonality preserving linear map between unital commutative
real C∗-algebras. As a consequence, we show that every orthogonality
preserving linear bijection between unital commutative real C∗-algebras is
continuous.
###### Key words and phrases:
Orthogonal form, real C∗-algebra, orthogonality preservers, disjointness
preserver, separating map
###### 2010 Mathematics Subject Classification:
Primary 46H40; 4J10, Secondary 47B33; 46L40; 46E15; 47B48.
Authors partially supported by the Spanish Ministry of Economy and
Competitiveness, D.G.I. project no. MTM2011-23843, and Junta de Andalucía
grants FQM0199 and FQM3737.
## 1\. Introduction and preliminaries
Elements $a$ and $b$ in a real or complex C∗-algebra, $A$, are said to be
_orthogonal_ (denoted by $a\perp b$) if $ab^{*}=b^{*}a=0$. A bounded bilinear
form $V:A\times A\to\mathbb{K}$ is called _orthogonal_ (resp., _orthogonal on
self-adjoint elements_) whenever $V(a,b^{*})=0$ for every $a\perp b$ in $A$
(resp., in the self-adjoint part of $A$). All the forms considered in this
paper are assumed to be continuous. Motivated by the seminal contributions by
K Ylinen [51] and R. Jajte and A. Paszkiewicz [29], S. Goldstein proved that
every orthogonal form $V$ on a (complex) C∗-algebra, $A,$ is of the form
$V(x,y)=\phi(xy)+\psi(xy)\ \ (x,y\in A),$
where $\phi$ and $\psi$ are two functionals in $A^{*}$ (cf. [21, Theorem
1.10]). A simplified proof of Goldstein’s theorem was published by U. Haagerup
and N.J. Laustsen in [24]. This characterisation has emerged as a very useful
tool in the study of bounded linear operators between C∗-algebras which are
orthogonality or disjointness preserving (see, for example, [11, 12]).
The first aim of this paper is to study orthogonal forms on the wider class of
real C∗-algebras. Little or nothing is known about the structure of an
orthogonal form $V$ on a real C∗-algebra. At first look, one is tempted to
consider the canonical complex bilinear extension of $V$ to a form on the
complexification, $A_{\mathbb{C}}=A\oplus iA$, of $A$ and, when the latter is
orthogonal, to apply Goldstein’s theorem. However, the complex bilinear
extension of $V$ to $A_{\mathbb{C}}\times A_{\mathbb{C}}$, need not be, in
general, orthogonal (see Example 2.7). The study of orthogonal forms on real
C∗-algebras requires a completely independent strategy; surprisingly the
resulting forms will enjoy a different structure to that established by S.
Goldstein in the complex setting.
In section 2 we establish some structure results for orthogonal forms on a
general real C∗-algebra, showing, among other properties, that every
orthogonal form on a real C∗-algebra extends to an orthogonal form on its
multiplier algebra (see Proposition 1.3). It is also proved that, for each
orthogonal and symmetric form $V$ on a real C∗-algebra, $A,$ there exists a
functional $\phi\in A^{*}$ satisfying $V(a,b)=\phi(ab+ba),$ for every $a,b\in
A$ with $a=a^{*}$, $b^{*}=b$ (cf. Proposition 1.5). In the real setting, the
skew-symmetric part of a real C∗-algebra, $A,$ is not determined by the self-
adjoint part of $A,$ so the information about the behavior of $V$ on the rest
of $A$ is very limited.
Section 3 contains one of the main results of the paper: the characterisation
of all orthogonal forms on a commutative real C∗-algebra. Concretely, we prove
that a form $V$ on a commutative real C∗-algebra $A$ is orthogonal if, and
only if, there exist functionals $\varphi_{1}$ and $\varphi_{2}$ in $A^{*}$
satisfying
$V(x,y)=\varphi_{1}(xy)+\varphi_{2}(xy^{*}),$
for every $x,y\in A$ (see Theorem 2.4). Among the consequences, it follows
that the complex bilinear extension of $V$ to the complexification of $A$ is
orthogonal if, and only if, we can take $\varphi_{2}=0$ in the above
representation.
We recall that a mapping $T:A\to B$ between real or complex C∗-algebras is
said to be _orthogonality or disjointness preserving_ (also called
_separating_) whenever $a\perp b$ in $A$ implies $T(a)\perp T(b)$ in $B.$ The
mapping $T$ is _bi-orthogonality preserving_ whenever the equivalence
$a\perp b\Leftrightarrow T(a)\perp T(b)$
holds for all $a,b$ in $A$. As noticed in [13], every bi-orthogonality
preserving linear surjection, $T:A\to B$ between two C∗-algebras is injective.
The study of orthogonality preserving operators between C∗-algebras started
with the work of W. Arendt [1] in the setting of unital abelian C∗-algebras.
Subsequent contributions by K. Jarosz [30] extended the study to the setting
of orthogonality preserving (not necessarily bounded) linear mappings between
abelian C∗-algebras. The first study on orthogonality preserving symmetric
(bounded) linear operators between general (complex) C∗-algebras is originally
due to M. Wolff (cf. [49]). Orthogonality preserving bounded linear maps
between C∗-algebras, JB∗-algebras and JB∗-triples were completely described in
[11] and [12].
The pioneer works of E. Beckenstein, L. Narici, and A.R. Todd in [8] and [9]
(see also [7]) were applied by K. Jarosz to prove that every orthogonality
preserving linear bijection between $C(K)$-spaces is (automatically)
continuous (see [30]). More recently, M. Burgos and the authors of this note
proved in [13] that every bi-orthogonality preserving linear surjection
between two von Neumann algebras (or between two compact C∗-algebras) is
automatically continuous (compare [40], [41] for recent additional
generalisations).
The main goal of section 4 is to describe the orthogonality preserving linear
mappings between unital commutative real C∗-algebras (see Theorem 3.2). As a
consequence, we shall prove that every orthogonality preserving linear
bijection between unital commutative real C∗-algebras is automatically
continuous. We shall exhibit some examples illustrating that the results in
the real setting are completely independent from those established for complex
C∗-algebras. We further give a characterisation of those linear mappings
between real forms of $C(K)$-spaces which are bi-orthogonality preserving.
### 1.1. Preliminary results
Let us now introduce some basic facts and definitions required later. A _real
C ∗-algebra_ is a real Banach *-algebra $A$ which satisfies the standard
C∗-identity, $\|a^{*}a\|=\|a\|^{2}$, and which also has the property that
$1+a^{\ast}a$ is invertible in the unitization of $A$ for every $a\in A$. It
is known that a real Banach *-algebra, $A,$ is a real C∗-algebra if, and only
if, it is isometrically *-isomorphic to a norm-closed real *-subalgebra of
bounded operators on a real Hilbert space (cf. [39, Corollary 5.2.11]).
Clearly, every (complex) C∗-algebra is a real C∗-algebra when scalar
multiplication is restricted to the real field. If $A$ is a real C∗-algebra
whose algebraic complexification is denoted by ${B}=A\oplus iA,$ then there
exists a C∗-norm on ${B}$ extending the norm of $A$. It is further known that
there exists an involutive conjugate-linear ∗-automorphism $\tau$ on ${B}$
such that $A={B}^{\tau}:=\\{x\in B:\tau(x)=x\\}$ (compare [39, Proposition
5.1.3] or [47, Lemma 4.1.13], and [22, Corollary 15.4]). The dual space of a
real or complex C∗-algebra $A$ will be denoted by $A^{*}$. Let
$\widetilde{\tau}:B^{*}\rightarrow B^{*}$ denote the map defined by
$\widetilde{\tau}(\phi)(b)=\overline{\phi(\tau(b))}\qquad(\phi\in B^{*},\,b\in
B).$
Then $\widetilde{\tau}$ is a conjugate-linear isometry of period 2 and the
mapping
$(B^{*})^{\widetilde{\tau}}\to A^{*}$ $\varphi\mapsto\varphi|_{A}$
is a surjective linear isometry. We shall identify
$(B^{*})^{\widetilde{\tau}}$ and $A^{*}$ without making any explicit mention.
When $A$ is a real or complex C∗-algebra, then $A_{sa}$ and $A_{skew}$ will
stand for the set of all self-adjoint and skew-symmetric elements in $A$,
respectively. We shall make use of standard notation in C∗-algebra theory.
Given Banach spaces $X$ and $Y$, $L(X,Y)$ will denote the space of all bounded
linear mappings from $X$ to $Y$. We shall write $L(X)$ for the space $L(X,X)$.
Throughout the paper the word “operator” (respectively, multilinear or
sesquilinear operator) will always mean bounded linear mapping (respectively
bounded multilinear or sesquilinear mapping). The dual space of a Banach space
$X$ is always denoted by $X^{*}$.
Let us recall that a series $\sum_{n}x_{n}$ in a Banach space is called
_weakly unconditionally Cauchy (w.u.C.)_ if there exists $C>0$ such that for
any finite subset $F\subset{\mathbb{N}}$ and $\varepsilon_{n}=\pm 1$ we have
$\displaystyle\left\|\sum_{n\in F}\varepsilon_{n}x_{n}\right\|\leq C$. A
(linear) operator $T:X\longrightarrow Y$ is _unconditionally converging_ if
for every w.u.C. series $\sum_{n}x_{n}$ in $X,$ the series $\sum_{n}T(x_{n})$
is unconditionally convergent in $Y$, that is, every subseries of
$\sum_{n}T(x_{n})$ is norm converging. It is known that $T:X\to Y$ is
unconditionally converging if, and only if, for every w.u.C. series
$\sum_{n}x_{n}$ in $X,$ we have $\|T(x_{n})\|\to 0$ (compare, for example,
[44, page 1257])
Let us also recall that a Banach space $X$ is said to have _Pełczyński’s_
_property_ (V) if, for every Banach space $Y$, every unconditionally
converging operator $T:X\to Y$ is weakly compact.
The proof of the following elementary lemma is left to the reader.
###### Lemma 1.1.
Let $X$ be a complex Banach space, $\tau:X\to X$ a conjugate-linear period-2
isometry. Then the real Banach space $X^{\tau}:=\\{x\in X:\tau(x)=x\\}$
satisfies property (V) whenever $X$ does. $\hfill\Box$
We shall require, for later use, some results on extensions of multilinear
operators. Let $X_{1},\dots,X_{n}$, and $X$ be Banach spaces,
$T:X_{1}\times\cdots\times X_{n}\to X$ a (continuous) $n$-linear operator, and
$\pi:\\{1,\dots,n\\}\to\\{1,\dots,n\\}$ a permutation. It is known that there
exists a unique $n$-linear extension $AB(T)_{\pi}:X_{1}^{**}\times\cdots\times
X_{n}^{**}\to X^{**}$ such that for every $z_{i}\in X_{i}^{**}$ and every net
$(x_{\alpha_{i}}^{i})\in X_{i}$ ($1\leq i\leq n$), converging to $z_{i}$ in
the weak* topology we have
$AB(T)_{\pi}(z_{1},\ldots,z_{n})=\mbox{weak*-}\lim_{\alpha_{\pi(1)}}\cdots\mbox{weak*-}\lim_{\alpha_{\pi(n)}}T(x_{\alpha_{1}}^{1},\ldots,x_{\alpha_{n}}^{n}).$
Moreover, $AB(T)_{\pi}$ is bounded and has the same norm as $T$. The
extensions $AB(T)_{\pi}$ coincide with those considered by Arens in [2, 3] and
by Aron and Berner for polynomials in [4]. The $n$-linear operators
$AB(T)_{\pi}$ are usually called the Arens or _Aron-Berner extensions_ of $T$.
Under some additional hypothesis, the Arens extension of a multilinear
operator also is separately weak∗ continuous. Indeed, if every operator from
$X_{i}$ to $X_{j}^{*}$ is weakly compact ($i\not=j)$ the Arens extensions of
$T$ defined above do not depend on the chosen permutation $\pi$ and they are
all separately weak∗ continuous (see [5], and Theorem 1 in [10]). In
particular, the above requirements always hold when every $X_{i}$ satisfies
Pelczynski’s property $(V)$ (in such case $X_{i}^{*}$ contains no copies of
$c_{0}$, therefore every operator from $X_{i}$ to $X_{j}^{*}$ is
unconditionally converging, and hence weakly compact by property $(V)$, see
[43]). When all the Arens extensions of $T$ coincide, the symbol
$AB(T)=T^{**}$ will denote any of them.
We should note at this point that every C∗-algebra satisfies property $(V)$
(cf. Corollary 6 in [46]). Since every real C∗-algebra is, in particular, a
real form of a (complex) C∗-algebra, it follows from Lemma 1.1 that every real
C∗-algebra satisifes property $(V)$. We therefore have:
###### Lemma 1.2.
Let $A_{1},\ldots,A_{k}$ be real C∗-algebras and let $T$ be a multilinear
continuous operator from $A_{1}\times\ldots\times A_{k}$ to a real Banach
space $X$. Then $T$ admits a unique Arens extension
$T^{**}:A_{1}^{**}\times\ldots\times A_{k}^{**}\to X^{**}$ which is separately
weak∗ continuous.$\hfill\Box$
Given a real or complex C∗-algebra, $A$, the _multiplier algebra_ of $A$,
$M(A)$, is the set of all elements $x\in A^{**}$ such that, for each element
$a\in A$, $xa$ and $ax$ both lie in $A$. We notice that $M(A)$ is a C∗-algebra
and contains the unit element of $A^{**}$. It should be recalled here that
$A=M(A)$ whenever $A$ is unital.
###### Proposition 1.3.
Let $A$ be a real C∗-algebra. Suppose that $V:A\times A\to\mathbb{R}$ is an
orthogonal bounded bilinear form. Then the continuous bilinear form
$\tilde{V}:M(A)\times M(A)\to\mathbb{R},\ \ \tilde{V}(a,b):=V^{**}(a,b)$
is orthogonal.
###### Proof.
Let $a$ and $b$ be two orthogonal elements in $M(A)$. Let $a^{[\frac{1}{3}]}$
(resp., $b^{[\frac{1}{3}]}$) denote the unique element $z$ in $M(A)$
satisfying $zz^{*}z=a$ (resp., $zz^{*}z=b$). We notice that
$a^{[\frac{1}{3}]}$ and $b^{[\frac{1}{3}]}$ are orthogonal, so, for each pair
$x,y$ in $A$, $a^{[\frac{1}{3}]}xa^{[\frac{1}{3}]}$ and
$b^{[\frac{1}{3}]}yb^{[\frac{1}{3}]}$ are orthogonal elements in $A$. Since
$V$ is orthogonal, we have
$V(a^{[\frac{1}{3}]}xa^{[\frac{1}{3}]},(b^{[\frac{1}{3}]})^{*}y(b^{[\frac{1}{3}]})^{*})=0$
for every $x,y\in A$.
Goldstine’s theorem (cf. Theorem V.4.2.5 in [18]) guarantees that the closed
unit ball of $A$ is weak*-dense in the closed unit ball of $A^{**}$. Therefore
we can pick two bounded nets $(x_{\lambda})$ and $(y_{\mu})$ in $A$,
converging in the weak∗ topology of $A^{**}$ to $(a^{[\frac{1}{3}]})^{*}$ and
$b^{[\frac{1}{3}]}$, respectively.
We have already mentioned that $V^{**}:A^{**}\times A^{**}\to\mathbb{R}$ is
separately weak∗ continuous. Since
$0=V(a^{[\frac{1}{3}]}x_{\lambda}a^{[\frac{1}{3}]},(b^{[\frac{1}{3}]})^{*}y_{\mu}(b^{[\frac{1}{3}]})^{*}),$
for every $\lambda$ and $\mu$, taking limits, first in $\lambda$ and
subsequently in $\mu$, we deduce that
$V^{**}(a^{[\frac{1}{3}]}(a^{[\frac{1}{3}]})^{*}a^{[\frac{1}{3}]},(b^{[\frac{1}{3}]})^{*}b^{[\frac{1}{3}]}(b^{[\frac{1}{3}]})^{*}))=\tilde{V}(a,b^{*})=0,$
which shows that $\tilde{V}$ is orthogonal. ∎
Since the multiplier algebra of a real or complex C∗-algebra always has a unit
element, Proposition 1.3 allows us to restrict our study on orthogonal
bilinear forms on a real C∗-algebra $A$ to the case in which $A$ is unital.
A _real von Neumann algebra_ is a real C∗-algebra which is also a dual Banach
space (cf. [28] or [39, §6.1]). Clearly, the self adjoint part of a real von
Neumann algebra is a JW-algebra in the terminology employed in [25], so every
self-adjoint element in a real von Neumann algebra $W$ can be approximated in
norm by a finite real linear combination of mutually orthogonal projections in
$W$ (cf. [25, Proposition 4.2.3]). We shall explore now the validity in the
real setting of some of the results established by S. Goldstein in [21].
###### Lemma 1.4.
Let $A$ be a real von Neumann algebra with unit $1$. Suppose that $V:A\times
A\to\mathbb{R}$ is a bounded bilinear form. The following are equivalent:
1. $(a)$
$V$ is orthogonal on $A_{sa}$;
2. $(b)$
$V(p,q)=0$, whenever $p$ and $q$ are two orthogonal projections in $A$;
3. $(c)$
$V(a,b)=V(ab,1)$ for every $a,b\in A_{sa}$ with $ab=ba$.
If any of the above statements holds and $V$ is symmetric, then defining
$\phi_{1}(x):=V(x,1)$ ($x\in A$), we have $V(a,b)=\phi_{1}(\frac{ab+ba}{2})$,
for every $a,b\in A_{sa}$.
###### Proof.
Applying the existence of spectral resolutions for self-adjoint elements in a
real von Neumann algebra, the argument given by S. Goldstein in [21,
Proposition 1.2] remains valid to prove the equivalence of $(a)$, $(b)$ and
$(c)$.
Suppose now that $V$ is symmetric. Let $a=\sum_{j=1}^{m}\lambda_{j}p_{j}$ be
an algebraic element in $A_{sa}$, where the $\lambda_{j}$’s belong to
$\mathbb{R}$ and $p_{1},\ldots,p_{m}$ are mutually orthogonal projections in
$A$. Since $V$ is orthogonal, for every projection $p\in A$, we have
$V(p,1)=V(p,1-p)+V(p,p)=V(p,p).$
Thus,
$V(a,a)=\sum_{j=1}^{m}\lambda_{j}^{2}V(p_{j},p_{j})=\sum_{j=1}^{m}\lambda_{j}^{2}V(p_{j},1)=V\left(\sum_{j=1}^{m}\lambda_{j}^{2}p_{j},1\right)=V(a^{2},1).$
The (norm) density of algebraic elements in $A_{sa}$ and the continuity of $V$
imply that $V(a,a)=V(a^{2},1),$ for every $a\in A_{sa}$. Finally, applying
that $V$ is symmetric we have
$V(a^{2},1)+V(b^{2},1)+V(ab+ba,1)=V((a+b)^{2},1)$
$=V(a+b,a+b)=V(a,a)+V(b,b)+2V(a,b),$
for every $a,b\in A_{sa}$, and hence $V(a,b)=V(\frac{ab+ba}{2},1),$ for all
$a,b\in A_{sa}$. ∎
The above result holds for every monotone $\sigma$-complete unital real
C∗-algebra $A$ (that is, each upper bounded, monotone increasing sequence of
selfadjoint elements of $A$ has a least upper bound).
Surprisingly, the final conclusion of the above Lemma can be established for
unital real C∗-algebras with independent basic techniques.
###### Proposition 1.5.
Let $A$ be a unital real C∗-algebra with unit $1$. Suppose that $V:A\times
A\to\mathbb{R}$ is an orthogonal, symmetric, bounded, bilinear form. Then
defining $\phi_{1}(x):=V(x,1)$ ($x\in A$), we have
$V(a,b)=\phi_{1}(\frac{ab+ba}{2})$, for every $a,b\in A_{sa}$.
###### Proof.
Let $a$ be a selfadjoint element in $A$. The real C∗-subalgebra, $C$, of $A$
generated by $1$ and $a$ is isometrically isomorphic to the space
$C(K,\mathbb{R})$ of all real-valued continuous functions on a compact
Hausdorff space $K$. The restriction of $V$ to $C\times C$ is orthogonal,
therefore the mapping $x\mapsto V(x,x)$ is a $2$-homogeneous orthogonally
additive polynomial on $C$. The main result in [45] implies the existence of a
functional $\varphi_{a}\in C^{*}$ such that $V(x,x)=\varphi_{a}(x^{2})$, for
every $x\in C$. It is clear that $\varphi_{a}(x)=V(x,1)$ for every $x\in C$.
In particular
$V(a,a)=\varphi_{a}(a^{2})=V(a^{2},1).$
The argument given at the end of the proof of Lemma 1.4 gives the desired
statement. ∎
The above proposition shows that we can control the form of a symmetric
orthogonal form on the self adjoint part of a (unital) real C∗-algebra. The
form on the skew-symmetric part remains out of control for the moment.
## 2\. Orthogonal forms on abelian real C∗-algebras
Throughout this section, $A$ will denote a unital, abelian, real C∗-algebra
whose complexification will be denoted by $B$. It is clear that $B$ is a
unital, abelian C∗-algebra. It is known that there exists a period-2
conjugate-linear ∗-automorphism $\tau:B\to B$ such that $A=B^{\tau}:=\\{x\in
B:\tau(x)=x\\}$ (cf. [47, 4.1.13] and [22, 15.4] or [39, §5.2]).
By the commutative Gelfand theory, there exists a compact Hausdorff space $K$
such that $B$ is C∗-isomorphic to the C∗-algebra $C(K)$ of all complex valued
continuous functions on $K$. The Banach-Stone Theorem implies the existence of
a homeomorphism $\sigma:K\to K$ such that $\sigma^{2}(t)=t$, and
$\tau(a)(t)=\overline{a(\sigma(t))},$
for all $t\in K$, $a\in C(K)$. Real function algebras of the form
$C(K)^{\tau}$ have been studied by its own right and are interesting in some
other settings (cf. [37]).
Henceforth, the symbol $\mathfrak{B}$ will stand for the $\sigma$-algebra of
all Borel subsets of $K$, $S(K)$ will denote the space of
$\mathfrak{B}$-simple scalar functions defined on $K$, while the _Borel
algebra over $K$_, $B(K)$, is defined as the completion of $S(K)$ under the
supremum norm. It is known that $B=C(K)\subset B(K)\subset C(K)^{**}.$ The
mapping $\tau^{**}:C(K)^{**}\to C(K)^{**}$ is a period-2 conjugate-linear
∗-automorphism on $B^{**}=C(K)^{**}.$ It is easy to see that
$\tau^{**}(B(K))=B(K),$ and hence $\tau^{**}|_{B(K)}:B(K)\to B(K)$ defines a
period-2 conjugate-linear ∗-automorphism on $B(K)$. By an abuse of notation,
the symbol $\tau$ will denote $\tau$, $\tau^{**}$ and $\tau^{**}|_{B(K)}$
indistinctly. It is clear that, for each Borel set $B\in\mathfrak{B},$
$\tau(\chi_{{}_{B}})=\chi_{{}_{\sigma(B)}}$.
Let $a$ be an element in $B(K)$. For each $\varepsilon>0,$ there exist complex
numbers $\lambda_{1},\ldots,\lambda_{r}$ and disjoint Borel sets
$B_{1},\ldots,B_{r}$ such that
$\displaystyle{\left\|a-\sum_{k=1}^{r}\lambda_{k}\chi_{{}_{B_{k}}}\right\|<{\varepsilon}.}$
When $a\in A$ is $\tau$-symmetric (i.e. $\tau(a)=a$) then, since
$a=\frac{1}{2}(a+\tau(a)),$ we have
$\left\|a-\frac{1}{2}\sum_{k=1}^{r}\lambda_{k}\chi_{{}_{B_{k}}}+\overline{\lambda_{k}}\chi_{{}_{\sigma(B_{k})}}\right\|\leq\frac{1}{2}\left\|a-\sum_{k=1}^{r}\lambda_{k}\chi_{{}_{B_{k}}}\right\|+\frac{1}{2}\left\|a-\sum_{k=1}^{r}\overline{\lambda_{k}}\chi_{{}_{\sigma(B_{k})}}\right\|$
$\leq\frac{1}{2}\left\|a-\sum_{k=1}^{r}\lambda_{k}\chi_{{}_{B_{k}}}\right\|+\frac{1}{2}\left\|\tau\left(a-\sum_{k=1}^{r}\lambda_{k}\chi_{{}_{B_{k}}}\right)\right\|<\varepsilon.$
Consequently, every element in $B(K)^{\tau}$ can be approximated in norm by
finite linear combinations of the form
$\sum_{k}\alpha_{k}\chi_{{}_{B_{k}}}+\overline{\alpha_{k}}\chi_{{}_{\sigma(B_{k})}},$
where $\alpha_{1},\ldots,\alpha_{n}$ are complex numbers and
$B_{1},\ldots,B_{n}$ are mutually disjoint Borel sets. Having in mind that,
for each Borel set $B\in\mathfrak{B}$ and each $\alpha\in\mathbb{C},$
$\left(\alpha\chi_{{}_{B}}+\overline{\alpha}\chi_{{}_{\sigma(B)}}\right)^{*}=\overline{\alpha}\chi_{{}_{B}}+{\alpha}\chi_{{}_{\sigma(B)}},$
we have
$\left(\alpha\chi_{{}_{B}}+\overline{\alpha}\chi_{{}_{\sigma(B)}}\right)+\left(\alpha\chi_{{}_{B}}+\overline{\alpha}\chi_{{}_{\sigma(B)}}\right)^{*}=2\Re\hbox{e}(\alpha)\left(2\chi_{{}_{\sigma(B)\cap
B}}+\chi_{{}_{\sigma(B)\backslash B}}+\chi_{{}_{B\backslash\sigma(B)}}\right)$
$=2\Re\hbox{e}(\alpha)\left(2\chi_{{}_{\sigma(B)\cap
B}}+\chi_{{}_{(\sigma(B)\backslash B)\cup\sigma(\sigma(B)\backslash
B)}}\right),$
and
$\left(\alpha\chi_{{}_{B}}+\overline{\alpha}\chi_{{}_{\sigma(B)}}\right)-\left(\alpha\chi_{{}_{B}}+\overline{\alpha}\chi_{{}_{\sigma(B)}}\right)^{*}=2i\Im\hbox{m}(\alpha)\left(\chi_{{}_{B\backslash\sigma(B)}}-\chi_{{}_{\sigma(B)\backslash
B}}\right).$
Suppose now that $a\in B(K)^{\tau}$ is *-symmetric (i.e. $a^{*}=a$). It
follows from the above that $a$ can be approximated in norm by linear
combinations of the form
$\displaystyle{\sum_{k=1}^{r}\alpha_{k}\chi_{{}_{E_{k}}},}$ where
$\alpha_{k}\in\mathbb{R}$ and $E_{1},\ldots,E_{r}$ are mutually disjoint Borel
subsets of $K$ with $\sigma(E_{i})=E_{i}$. Let $b$ be an element in
$B(K)^{\tau}$ satisfying $b^{*}=-b$. Similar arguments to those given for
*-symmetric elements, allow us to show that $b$ can be approximated in norm by
finite linear combinations of the form $\displaystyle{\sum_{k=1}^{r}i\
\alpha_{k}(\chi_{{}_{E_{k}}}-\chi_{{}_{\sigma(E_{k})}}),}$ where
$\alpha_{k}\in\mathbb{R}$ and $E_{1},\ldots,E_{r}$ are mutually disjoint Borel
subsets of $K$ with $\sigma(E_{i})\cap E_{i}=\emptyset$.
###### Lemma 2.1.
Let $A$ be a unital, abelian, real C∗-algebra whose complexification is
denoted by $B=C(K)$, for a suitable compact Hausdorff space $K$. Let
$\tau:B\to B$ be a period-2 conjugate-linear ∗-automorphism satisfying
$A=B^{\tau}$ and $\tau(a)(t)=\overline{a(\sigma(t))},$ for all $t\in K$, $a\in
C(K)$, where $\sigma:K\to K$ is a period-2 homeomorphism. Then the set
$N=\\{t\in K:\sigma(t)\neq t\\}$ is an open subset of $K$, $F=\\{t\in
K:\sigma(t)=t\\}$ is a closed subset of $K$ and there exists an open subset
$\mathcal{O}\subset K$ maximal with respect to the property
$\mathcal{O}\cap\sigma(\mathcal{O})=\emptyset.$
###### Proof.
That $F$ is closed follows easily from the continuity of $\sigma,$ and
consequently, $N=K/F$ is open.
Let $\mathcal{F}$ be the family of all open subsets $O\subseteq K$ such that
$O\cap\sigma(O)=\emptyset$ ordered by inclusion. Let
$S=\\{O_{\lambda}\\}_{\lambda}$ be a totally ordered subset of $\mathcal{F}.$
We shall see that $O=\bigcup_{\lambda}O_{\lambda}$ is an open set which also
lies in $\mathcal{F},$ that is, $O\cap\sigma(O)=\emptyset.$
Let us suppose, on the contrary, that there exists $t\in
O\cap\sigma(O)\neq\emptyset.$ Then there exist $\lambda,\beta$ such that $t\in
O_{\lambda}$ and $t\in\sigma(O_{\beta}).$ Since $S$ is totally ordered,
$O_{\lambda}\subseteq O_{\beta}$ or $O_{\beta}\subseteq O_{\lambda}.$ We shall
assume that $O_{\lambda}\subseteq O_{\beta}.$ Then
$\sigma(O_{\lambda})\subseteq\sigma(O_{\beta})$ and $t$ lies in
$O_{\beta}\cap\sigma(O_{\beta})=\emptyset,$ which is a contradiction. Finally,
Zorn’s Lemma gives the existence of a maximal element $\mathcal{O}$ in
$\mathcal{F}.$ ∎
It should be noticed here that, in Lemma 2.1,
$\mathcal{O}\cup\sigma(\mathcal{O})=N$, an equality which follows from the
maximality of $\mathcal{O}$.
Our next lemma analyses the “spectral resolution” of a *-skew-symmetric
element in $B(K)^{\tau}.$
###### Lemma 2.2.
In the notation of Lemma 2.1, let $B(A)=B(K)^{\tau}$, let $a\in
B(K)^{\tau}_{sa},$ and let $b$ be an element in $B(A)_{skew}$. Then the
following statements hold:
1. $a)$
$b|F=0$;
2. $b)$
For each $\varepsilon>0$, there exist mutually disjoint Borel sets
$B_{1},\ldots,B_{m}\subset\mathcal{O}$ and real numbers
$\lambda_{1},\ldots,\lambda_{m}$ satisfying
$\displaystyle{\left\|b-\sum_{j=1}^{m}i\
\lambda_{j}(\chi_{{}_{B_{j}}}-\chi_{{}_{\sigma(B_{j})}})\right\|<\varepsilon;}$
3. $c)$
For each $\varepsilon>0$, there exist mutually disjoint Borel sets
$C_{1},\ldots,C_{m}\subset K$ and real numbers $\mu_{1},\ldots,\mu_{m}$
satisfying $\sigma(C_{j})=C_{j},$ and
$\displaystyle{\left\|a-\sum_{j=1}^{m}\mu_{j}\chi_{{}_{C_{j}}}\right\|<\varepsilon.}$
###### Proof.
$a)$ Since $b^{*}=-b$, we have $Re(b(t))=0,\forall t\in K.$ Now, let $t\in F$,
applying $\sigma(t)=t$ and $\tau(b)=b$ we get
$\overline{b(t)}=\overline{b(\sigma(t))}=b(t),$ and hence
$\Im\hbox{m}(b(t))=0.$
Statements $b)$ and $c)$ follow from the comments prior to Lemma 2.1 and the
maximality of $\mathcal{O}$ in that Lemma. ∎
It is clear that in a commutative real (or complex) C∗-algebra, $A,$ two
elements $a,b$ are orthogonal if and only if they have zero-product, that is,
$ab=0.$ Therefore, $V(a,b^{*})=0=V(a,b)$ whenever $V:A\times A\to{\mathbb{R}}$
is an orthogonal bilinear form on an abelian real C∗-algebra and $a,b$ are two
orthogonal elements in $A$. We shall make use of this property without an
explicit mention.
We shall keep the notation of Lemma 2.1 throughout the section. Henceforth,
for each $C\subseteq\mathcal{O}$ we shall write $u_{{}_{C}}=i\
(\chi_{{}_{C}}-\chi_{{}_{\sigma(C)}}).$ The symbol $u_{{}_{0}}$ will stand for
the element $u_{{}_{\mathcal{O}}}.$ It is easy to check
$1=\chi_{{}_{F}}+u_{{}_{0}}u_{{}_{0}}^{*},$ where $1$ is the unit element in
$B(K)^{\tau}.$ By Lemma 2.2 $a)$, for each $b\in B(K)^{\tau}_{skew}$ we have
$b\perp\chi_{{}_{F}},$ and so $b=bu_{{}_{0}}u_{{}_{0}}^{*}.$
###### Proposition 2.3.
Let $K$ be a compact Hausdorff space, $\tau$ a period-2 conjugate-linear
isometric ∗-homomorphism on $C(K)$, $A=C(K)^{\tau}$, and $V:A\times
A\to\mathbb{R}$ be an orthogonal bounded bilinear form whose Arens extension
is denoted by $V^{**}:A^{**}\times A^{**}\to\mathbb{R}$. Let $\sigma:K\to K$
be a period-2 homeomorphism satisfying $\tau(a)(t)=\overline{a(\sigma(t))},$
for all $t\in K$, $a\in C(K)$. Then the following assertions hold for all
Borel subsets $D,B,C$ of $K$ with $\sigma(B)\cap B=\sigma(C)\cap C=\emptyset$
and $\sigma(D)=D$:
1. $a)$
$V(\chi_{{}_{D}},u_{{}_{B}})=V(u_{{}_{B}},\chi_{D})=0,$ whenever $D\cap
B=\emptyset$;
2. $b)$
$V(u_{{}_{B}},u_{{}_{C}})=0,$ whenever $B\cap C=\emptyset;$
3. $c)$
$V((u_{{}_{0}}u_{{}_{0}}^{*}-u_{{}_{C}}u^{*}_{{}_{C}})u_{{}_{B}},u_{{}_{C}})=V(u_{{}_{C}},(u_{{}_{0}}u_{{}_{0}}^{*}-u_{{}_{C}}u^{*}_{{}_{C}})u_{{}_{B}})=0.$
###### Proof.
By an abuse of notation, we write $V$ for $V$ and $V^{**}$.
Let $K_{1},K_{2}$ be compact subsets of $K$ such that $K_{1},K_{2}$ and
$\sigma(K_{2})$ are mutually disjoint. By regularity and Urysohn’s Lemma there
exist nets $(f_{\lambda})_{{}_{\lambda}}$, $(g_{\gamma})_{{}_{\gamma}}$ in
$C(K)^{+}$ such that $\chi_{{}_{K_{1}}}\leq
f_{\lambda}\leq\chi_{{}_{K\backslash(K_{2}\cup\sigma(K_{2}))}},$
$\chi_{{}_{K_{2}}}\leq
g_{\gamma}\leq\chi_{{}_{K\backslash(K_{1}\cup\sigma(K_{1})\cup\sigma(K_{2}))}},$
$(f_{\lambda})_{{}_{\lambda}}$ (respectively, $(g_{\gamma})_{{}_{\gamma}}$)
converges to $\chi_{{}_{K_{1}}}$ (resp., to $\chi_{{}_{K_{2}}}$) in the weak∗
topology of $C(K)^{**}$.
The nets $\widetilde{f}_{\lambda}=\frac{1}{2}(f_{\lambda}+\tau(f_{\lambda}))$
and $\widetilde{g}_{\gamma}=i(g_{\gamma}-\tau(g_{\gamma}))$ lie in
$C(K)^{\tau}$ and converge in the weak∗ topology of $C(K)^{**}$ to
$\frac{1}{2}(\chi_{{}_{K_{1}}}+\chi_{{}_{\sigma(K_{1})}})$ and
$u_{{}_{K_{2}}},$ respectively. It is also clear that $f_{\lambda}\perp
g_{\gamma},$ $\tau(f_{\lambda})\perp g_{\gamma},$ and hence
$\widetilde{f}_{\lambda}\perp\widetilde{g}_{\gamma},$ for every
$\lambda,\gamma.$
By the separate weak∗ continuity of $V^{**}\equiv V$ we have
(1)
$V\left(\frac{1}{2}(\chi_{{}_{K_{1}}}+\chi_{{}_{\sigma(K_{1})}}),u_{{}_{K_{2}}}\right)=w^{*}-\lim_{\lambda}\left(w^{*}-\lim_{\gamma}V\left(\widetilde{f_{\lambda}},\widetilde{g}_{\gamma}\right)\right)=0,$
and
$V\left(u_{{}_{K_{2}}},\frac{1}{2}(\chi_{{}_{K_{1}}}+\chi_{{}_{\sigma(K_{1})}})\right)=0.$
We can similarly prove that
(2) $V\left(u_{{}_{K_{1}}},u_{{}_{K_{2}}}\right)=0,$
whenever $K_{1}$ and $K_{2}$ are two compact subsets of $K$ such that
$K_{1},K_{2},\sigma(K_{1})$ and $\sigma(K_{2})$ are pairwise disjoint.
$a)$ Let now $D,B$ be two disjoint Borel subsets of $K$ such that
$\sigma(D)=D$ and $B\subseteq\mathcal{O}.$ By inner regularity there exist
nets of the form $(\chi_{{}_{K_{\lambda}^{{}^{D}}}})_{\lambda}$ and
$(\chi_{{}_{K_{\gamma}^{{}^{B}}}})_{\gamma}$ such that
$(\chi_{{}_{K_{\lambda}^{{}^{D}}}})_{\lambda}$ and
$(\chi_{{}_{K_{\gamma}^{{}^{B}}}})_{\gamma}$ converge in the weak∗ topology of
$C(K)^{**}$ to $\chi_{{}_{D}}$ and $\chi_{{}_{B}}$, respectively, where each
$K_{\lambda}^{{}^{D}}\subseteq D$ and each $K_{\gamma}^{{}^{B}}\subseteq B$ is
compact subset of $K$. By the assumptions made on $D$ and $B$ we have that
$K_{\lambda}^{{}^{D}}\cap
K_{\gamma}^{{}^{B}}=K_{\lambda}^{{}^{D}}\cap\sigma(K_{\gamma}^{{}^{B}})=\emptyset$
and $K_{\gamma}^{{}^{B}}\subseteq\mathcal{O}$ for all $\lambda$ and $\gamma.$
By (1) and the separate weak∗ continuity of $V$ we have
(3)
$V(\chi_{{}_{D}},u_{{}_{B}})=w^{*}-\lim_{\lambda}\left(w^{*}-\lim_{\gamma}V\left(\frac{\chi_{{}_{K_{\lambda}^{{}^{D}}}}+\chi_{{}_{\sigma(K_{\lambda}^{{}^{D}})}}}{2},u_{{}_{K_{\gamma}^{{}^{B}}}}\right)\right)=0,$
and
(4) $V(u_{{}_{B}},\chi_{{}_{{}_{D}}})=0.$
A similar argument, but replacing (1) with (2), applies to prove $b)$.
To prove the last statement, we observe that
$(u_{{}_{0}}u_{{}_{0}}^{*}-u_{c}u_{c}^{*})u_{{}_{B}}=(\chi_{{}_{\mathcal{O}}}+\chi_{{}_{\sigma(\mathcal{O}})}-\chi_{{}_{C}}-\chi_{{}_{\sigma(C)}})u_{{}_{B}}=(\chi_{{}_{\mathcal{O}\setminus
C}}+\chi_{{}_{\sigma(\mathcal{O}\setminus
C)}})u_{{}_{B}}=u_{{}_{(\mathcal{O}\setminus C)\cap B}},$
and hence the statement $c)$ follows from $b)$. ∎
We can now establish the description of all orthogonal forms on a commutative
real C∗-algebra.
###### Theorem 2.4.
Let $V:A\times A\to\mathbb{R}$ be a continuous orthogonal form on a
commutative real C∗-algebra, then there exist $\varphi_{1},$ and $\varphi_{2}$
in $A^{*}$ satisfying
$V(x,y)=\varphi_{1}(xy)+\varphi_{2}(xy^{*}),$
for every $x,y\in A.$
###### Proof.
We may assume, without loss of generality, that $A$ is unital (compare
Proposition 1.3). Let $B$ denote the complexification of $A$. In this case $B$
identifies with $C(K)$ for a suitable compact Hausdorff space $K$ and
$A=C(K)^{\tau}$, where $\tau$ is a conjugate-linear period-2 *-homomorphism on
$C(K)$. We shall follow the notation employed in the rest of this section.
The form $V:A\times A\to\mathbb{R}$ extends to a continuous form
$V^{**}:A^{**}\times A^{**}\to\mathbb{R}$ which is separately weak∗ continuous
(cf. Lemma 1.2). The restriction $V^{**}|_{B(K)^{\tau}\times
B(K)^{\tau}}:B(K)^{\tau}\times B(K)^{\tau}\to\mathbb{R}$ also is a continuous
extension of $V$. We shall prove the statement for $V^{**}|_{B(K)^{\tau}\times
B(K)^{\tau}}.$ Henceforth, the symbol $V$ will stand for $V$, $V^{**}$ and
$V^{**}|_{B(K)^{\tau}\times B(K)^{\tau}}$ indistinctly.
Let us first take two self-adjoint elements $a_{1},a_{2}$ in $B(K)^{\tau}$. By
Proposition 1.5,
(5) $V(a_{1},a_{2})=V(a_{1}a_{2},1).$
To deal with the skew-symmetric part, let $D,B,C$ be Borel subsets of $K$
with, $D=\sigma(D)$ and $B,C\subseteq\mathcal{O}.$ From Proposition 2.3 $a)$,
we have
(6)
$V(\chi_{{}_{D}},u_{{}_{B}})=V(\chi_{{}_{D}},u_{{}_{B}}(1-\chi_{{}_{D}}+\chi_{{}_{D}}))=V(\chi_{{}_{D}},u_{{}_{B\cap(K\backslash
D)}})+V(\chi_{{}_{D}},u_{{}_{B}}\chi_{{}_{D}})$
$=V(\chi_{{}_{D}}-1+1,u_{{}_{B}}\chi_{{}_{D}})=V(-\chi_{{}_{(K\backslash
D)}}+1,u_{{}_{(B\cap D)}})=V(1,u_{{}_{B}}\chi_{{}_{D}}).$
Similarly,
(7) $V(u_{{}_{B}},\chi_{{}_{D}})=V(u_{{}_{B}}\chi_{{}_{D}},1).$
Now, Proposition 2.3 $b)$ and $c)$, repeatedly applied give:
$V(u_{{}_{B}},u_{{}_{C}})=V(u_{{}_{B}}(\chi_{{}_{F}}+u_{{}_{0}}u_{{}_{0}}^{*}),u_{{}_{C}})=V(u_{{}_{B}}u_{{}_{0}}u_{{}_{0}}^{*},u_{{}_{C}})$
$=V(u_{{}_{B}}(u_{{}_{0}}u_{{}_{0}}^{*}+u_{{}_{C}}u^{*}_{{}_{C}}-u_{{}_{C}}u^{*}_{{}_{C}}),u_{{}_{C}})=V(u_{{}_{B}}u_{{}_{C}}u^{*}_{{}_{C}},u_{{}_{C}})$
$=V(u_{{}_{B}}u_{{}_{C}}u^{*}_{{}_{C}},u_{{}_{C}}-u_{{}_{0}}+u_{{}_{0}})=V(u_{{}_{(B\cap
C)}},-u_{{}_{(\mathcal{O}\backslash C)}}+u_{{}_{0}})=V(u_{{}_{(B\cap
C)}},u_{{}_{0}})$
$=V(u_{{}_{B}}u_{{}_{C}}(u^{*}_{{}_{C}}-u_{{}_{0}}^{*}+u_{{}_{0}}^{*}),u_{{}_{0}})=V(u_{{}_{B}}u_{{}_{C}}u_{{}_{0}}^{*},u_{{}_{0}}).$
Thus, we have
(8)
$V(u_{{}_{B}},u_{{}_{C}})=V(u_{{}_{B}}u_{{}_{C}}u_{{}_{0}}^{*},u_{{}_{0}}),$
and similarly
(9)
$V(u_{{}_{B}},u_{{}_{C}})=V(u_{{}_{0}},u_{{}_{B}}u_{{}_{C}}u_{{}_{0}}^{*}).$
Let $\displaystyle{a_{l}=\sum_{j=1}^{m_{l}}\mu_{l,j}\chi_{{}_{D^{l}_{j}}},}$
$\displaystyle{b_{l}=\sum_{k=1}^{p_{l}}\lambda_{l,k}u_{{}_{B^{l}_{k}}}}$
($l\in\\{1,2\\}$) be two simple elements in $B(K)^{\tau}_{sa}$ and
$B(K)^{\tau}_{skew}$, respectively, where
$\lambda_{l,k},\mu_{l,j}\in\mathbb{R},$ for each $l\in\\{1,2\\},$
$\\{D^{l}_{1},\ldots,D^{l}_{m_{l}}\\}$ and
$\\{B^{l}_{1},\ldots,B^{l}_{p_{l}}\\}$ are families of mutually disjoint Borel
subsets of $K$ with $\sigma(D^{l}_{j})=D^{l}_{j}$ and
$B^{l}_{i}\subseteq\mathcal{O}.$ By $(\ref{eq theorem -1})$, $(\ref{eq sym
against anti})$, $(\ref{eq anti against sym})$, and $(\ref{eq anti anti 1}),$
we have
$V(a_{1}+b_{1},a_{2}+b_{2})=V(a_{1}a_{2},1)+\sum_{j=1}^{m_{1}}\sum_{k=1}^{p_{2}}\mu_{1,j}\lambda_{2,k}V\left(\chi_{{}_{D^{1}_{j}}},u_{{}_{B^{2}_{k}}}\right)$
$+\sum_{k=1}^{p_{1}}\sum_{j=1}^{m_{2}}\mu_{2,j}\lambda_{1,k}V\left(u_{{}_{B^{1}_{k}}},\chi_{{}_{D^{2}_{j}}}\right)+\sum_{k=1}^{p_{1}}\sum_{k=1}^{p_{2}}\lambda_{2,k}\lambda_{1,k}V\left(u_{{}_{B^{1}_{k}}},u_{{}_{B^{2}_{k}}}\right)$
$=V(a_{1}a_{2},1)+\sum_{j=1}^{m_{1}}\sum_{k=1}^{p_{2}}\mu_{1,j}\lambda_{2,k}V\left(1,\chi_{{}_{D^{1}_{j}}}u_{{}_{B^{2}_{k}}}\right)$
$+\sum_{k=1}^{p_{1}}\sum_{j=1}^{m_{2}}\mu_{2,j}\lambda_{1,k}V\left(u_{{}_{B^{1}_{k}}}\chi_{{}_{D^{2}_{j}}},1\right)+\sum_{k=1}^{p_{1}}\sum_{k=1}^{p_{2}}\lambda_{2,k}\lambda_{1,k}V\left(u_{{}_{B^{1}_{k}}}u_{{}_{B^{2}_{k}}}u_{{}_{0}}^{*},u_{{}_{0}}\right)$
$=V(a_{1}a_{2},1)+V\left(1,a_{1}b_{2}\right)+V\left(b_{1}a_{2},1\right)+V\left(b_{1}b_{2}u_{{}_{0}}^{*},u_{{}_{0}}\right)$
$=\psi_{1}(a_{1}a_{2})+\psi_{2}\left(a_{1}b_{2}\right)+\psi_{1}\left(b_{1}a_{2}\right)+\psi_{4}\left(b_{1}b_{2}\right),$
where $\psi_{1},\psi_{2},$ and $\psi_{4}$ are the functionals in $A^{*}$
defined by $\psi_{1}(x)=V(x,1),$ $\psi_{2}(x)=V(1,x),$ and
$\psi_{4}(x)=V(xu_{{}_{0}}^{*},u_{{}_{0}}),$ respectively. Since, by
Proposition 2.2, simple elements of the above form are norm-dense in
$B(K)^{\tau}_{sa}$ and $B(K)^{\tau}_{skew}$, respectively, and $V$ is
continuous, we deduce that
$V(a_{1}+b_{1},a_{2}+b_{2})=\psi_{1}(a_{1}a_{2})+\psi_{2}\left(a_{1}b_{2}\right)+\psi_{1}\left(b_{1}a_{2}\right)+\psi_{4}\left(b_{1}b_{2}\right),$
for every $a_{1},a_{2}\in B(K)^{\tau}_{sa}$, $b_{1},b_{2}\in
B(K)^{\tau}_{skew}.$
Now, taking $\phi_{1}=\frac{1}{4}(2\psi_{1}+\psi_{2}+\psi_{4}),$
$\phi_{2}=\frac{1}{4}(2\psi_{1}-\psi_{2}-\psi_{4}),$
$\phi_{3}=\frac{1}{4}(\psi_{2}-\psi_{4}),$ and
$\phi_{4}=\frac{1}{4}(\psi_{4}-\psi_{2}),$ we get
$V(a_{1}+b_{1},a_{2}+b_{2})=\phi_{1}((a_{1}+b_{1})(a_{2}+b_{2}))+\phi_{2}\left((a_{1}+b_{1})(a_{2}+b_{2})^{*}\right)$
$+\phi_{3}\left((a_{1}+b_{1})^{*}(a_{2}+b_{2})\right)+\phi_{4}\left((a_{1}+b_{1})^{*}(a_{2}+b_{2})^{*}\right),$
for every $a_{1},a_{2}\in B(K)^{\tau}_{sa}$, $b_{1},b_{2}\in
B(K)^{\tau}_{skew}.$
Finally, defining $\varphi_{1}(x)=\phi_{1}(x)+\phi_{4}(x^{*})$ and
$\varphi_{2}(x)=\phi_{2}(x)+\phi_{3}(x^{*})$ $(x\in A)$, we get the desired
statement. ∎
###### Remark 2.5.
The functionals $\varphi_{1}$ and $\varphi_{2}$ appearing in Theorem 2.4 need
not be unique. For example, let $(\varphi_{1},\varphi_{2})$ and
$(\phi_{1},\phi_{2})$ be two couples of elements in the dual of a commutative
real C∗-algebra $A$. It is not hard to check that
$\varphi_{1}(xy)+\varphi_{2}(xy^{*})=\phi_{1}(xy)+\phi_{2}(xy^{*}),$
for every $x,y\in A$ if, and only if,
$\varphi_{1}+\varphi_{2}=\phi_{1}+\phi_{2}$,
$(\varphi_{1}-\varphi_{2})(z)=(\phi_{1}-\phi_{2})(z)$ and
$(\varphi_{1}-\varphi_{2})(zw)=(\phi_{1}-\phi_{2})(zw),$ for every $z,w\in
A_{skew}.$ These conditions are not enough to guarantee that
$\phi_{i}=\varphi_{i}$. Take, for example,
$A=\mathbb{R}\oplus^{\infty}\mathbb{C}_{\mathbb{R}},$
$\phi_{1}(a,b)=a+\Re\hbox{e}(b)+\Im\hbox{m}(b),$ $\phi_{2}(a,b)=0,$
$\varphi_{1}(a,b)=\frac{a}{2}+\Re\hbox{e}(b)+\Im\hbox{m}(b),$ and
$\varphi_{2}(a,b)=\frac{a}{2}.$
###### Corollary 2.6.
Let $V:A\times A\to\mathbb{R}$ be a continuous orthogonal form on a
commutative real C∗-algebra, then its (unique) Arens extension
$V^{**}:A^{**}\times A^{**}\to\mathbb{R}$ is an orthogonal form.$\hfill\Box$
Clearly, the statement of the above Theorem 2.4 doesn’t hold for bilinear
forms on a commutative (complex) C∗-algebra. The real version established in
this paper is completely independent to the result proved by K. Ylinen for
commutative complex C∗-algebras in [51] and [21]. It seems natural to ask
whether the real result follows from the complex one by a mere argument of
complexification. Our next example shows that the (canonical) extension of an
orthogonal form on a commutative real C∗-algebra need not be an orthogonal
form on the complexification.
###### Example 2.7.
Let $K=\\{t_{1},t_{2}\\}.$ We define $\sigma:K\rightarrow K$ by
$\sigma(t_{1})=t_{2}.$ Let $A=C(K)^{\tau}$ be the real C∗-algebra whose
complexification is $C(K)$ and let $V:A\times A\to\mathbb{R},$ be the
orthogonal form defined by
$V(x,y)=\phi_{{}_{t_{1}}}(xy^{*})=\Re\hbox{e}(x(t_{1})\overline{y(t_{1})})=\Re\hbox{e}(x(t_{1})y(t_{2})),$
where $\phi_{t_{1}}=\Re\hbox{e}(\delta_{{}_{t_{1}}}).$ In this case, the
canonical complex bilinear extension $\widetilde{V}:C(K)\times
C(K)\to\mathbb{C}$ is given by
$\widetilde{V}(x,y)=\phi_{{}_{t_{1}}}(x\tau(y)^{*})=x(t_{1})y(t_{2})$ $(x,y\in
C(K)).$ It is clear that $\chi_{{}_{t_{1}}}\perp\chi_{{}_{t_{2}}}$ in $C(K),$
however $\widetilde{V}(\chi_{{}_{t_{1}}},\chi_{{}_{t_{2}}})=1\neq 0,$ which
implies that $\widetilde{V}$ is not orthogonal.
The (complex) bilinear extension of an orthogonal form $V$ on a real
C∗-algebra to its complexification is orthogonal precisely when $V$ satisfies
the generic form of an orthogonal form on a (complex) C∗-algebra given by the
main result in [21].
###### Corollary 2.8.
Let $V:A\times A\to\mathbb{R}$ be a continuous orthogonal form on a
commutative real C∗-algebra, let $B$ denote the complexification of $A$ and
let $\widetilde{V}:B\times B\to\mathbb{R}$ be the (complex) bilinear extension
of $V$. Then the form $\widetilde{V}$ is orthogonal if, and only if, $V$
writes in the form $V(x,y)=\varphi_{1}(xy)$ $(x,y\in A),$ where $\varphi_{1}$
is a functional in $A^{*}$.
###### Proof.
Let $\tau$ be the period-2 ∗-automorphism on $B$ satisfying that $B^{\tau}=B$
and let $\widetilde{\tau}:B^{*}\to B^{*}$ be the involution defined by
$\widetilde{\tau}(\phi)(b)=\overline{\phi(\tau(b))}$.
Suppose $\widetilde{V}$ is orthogonal. By the main result in [21] (see also
[51]), there exists $\phi\in B^{*}$ satisfying $\widetilde{V}(x,y)=\phi(xy),$
for every $x,y\in B$. Since $\widetilde{V}$ is an extension of $V$, we get
$V(a,b)=\Re\hbox{e}\phi(ab)=\phi(ab),$ for every $a,b\in A$. In particular,
$\phi(a)\in\mathbb{R}$, for every $a\in A$ and hence
$\widetilde{\tau}(\phi)=\phi$ lies in $(B^{*})^{\widetilde{\tau}}\equiv
A^{*}$.
Let us assume that $V$ writes in the form $V(x,y)=\varphi_{1}(xy)$ $(x,y\in
A)$, where $\varphi_{1}$ is a functional in $A^{*}$. The functional
$\varphi_{1}$ can be regarded as an element in $B^{*}$ satisfying
$\widetilde{\tau}(\varphi_{1})=\varphi_{1}.$ It is easy to check that
$\widetilde{V}(x,y)=\varphi_{1}(xy),$ for every $x,y\in B$. ∎
## 3\. Orthogonality preservers between commutative real C∗-algebras
Throughout this section, $A_{1}=C(K_{1})^{\tau_{1}}$ and
$A_{2}=C(K_{2})^{\tau_{2}}$ will denote two unital commutative real
C∗-algebras, $K_{1}$ and $K_{2}$ will be two compact Hausdorff spaces and
$\tau_{i}$ will denote a conjugate-linear period-2 ∗-automorphism on
$C(K_{i})$ given by $\tau_{i}(f)(t)=\overline{f(\sigma_{i}(t))}$ ($t\in
K_{i}$, $f\in C(K_{i})$), where $\sigma_{i}:K_{i}\to K_{i}$ is a period-2
homeomorphism. We shall write $B_{1}=C(K_{1})$ and $B_{2}=C(K_{2})$ for the
corresponding complexifications of $A_{1}$ and $A_{2}$, respectively.
By the Banach-Stone theorem, every surjective isometry $T:C(K_{1})\to
C(K_{2})$ is a composition operator, that is, there exist a unitary element
$u$ in $C(K_{2})$ and a homeomorphism $\sigma:K_{2}\to K_{1}$ such that
$T(f)(t)=(uC_{\sigma})(f)(t):=u(t)\ f(\sigma(t))$ ($t\in K_{2}$, $f\in
C(K_{1})$). This result led to the study of the so-called Banach-Stone
theorems in different classes of Banach spaces containing $C(K)$-spaces, in
which their algebraic and geometric properties are mutually determined. That
is the case of general C∗-algebras (R. Kadison [31] and Paterson and Sinclair
[42]), JB- and JB∗-algebras (Wright and M. Youngson [50] and Isidro and A.
Rodríguez [27]), JB∗-triples (Kaup [33] and Dang, Friedman and Russo [15]),
real C∗-algebras (Grzesiak [23], Kulkarni and Arundhathi [36], Kulkarni and
Limaye [37] and Chu, Dang, Russo and Ventura [14]) and real JB∗-triples
(Isidro, Kaup and Rodrguez [26], Kaup [34] and Fernndez-Polo, Martnez and
Peralta [19]). In what concerns us, we highlight that any surjective linear
isometry $T:C(K_{1})^{\tau_{1}}\to C(K_{2})^{\tau_{2}}$ is a composition
operator given by a homeomorphism $\phi:K_{2}\to K_{1}$ which satisfies
$\sigma_{1}\circ\phi=\phi\circ\sigma_{2}$ (cf. [23] or [36] or [37, Corollary
5.2.4]).
The class of orthogonality preserving (continuous) operators between
$C(K)$-spaces is strictly bigger than the class of surjective isometries.
Actually, a bounded linear operator $T:C(K_{1})\to C(K_{2})$ is orthogonality
preserving (equivalently, disjointness preserving) if, and only if, there
exist $u$ in $C(K_{2})$ and a mapping $\varphi:K_{2}\to K_{1}$ which is
continuous on $\\{t\in K_{2}:u(t)\neq 0\\}$ such that
$T(f)(t)=(uC_{\varphi})(f)(t)=u(t)\ f(\varphi(t))$ (compare [1, Example
2.2.1]).
Developing ideas given by E. Beckenstein, L. Narici, and A.R. Todd in [8] and
[9] (see also [7]), K. Jarosz showed, in [30], that the above hypothesis of
$T$ being continuous can be, in some sense, relaxed. More concretely, for
every orthogonality preserving linear mapping $T:C(K_{1})\to C(K_{2})$, there
exists a disjoint decomposition $K_{2}=S_{1}\cup S_{2}\cup S_{3}$ (with
$S_{2}$ open, $S_{3}$ closed), and a continuous mapping $\varphi$ from
$S_{1}\cup S_{2}$ into $K_{1}$ such that $T(f)(s)=\chi(s)f(\varphi(s))$ for
all $s\in S_{1}$ (where $\chi$ is a continuous, bounded, non-vanishing,
scalar-valued function on $S_{1}$), $T(f)(s)=0$ for all $s\in S_{3}$,
$\varphi(S_{2})$ is finite and, for each $s\in S_{2}$, the mapping $f\mapsto
T(f)(s)$ is not continuous. As a consequence, every orthogonality preserving
linear bijection between $C(K)$-spaces is (automatically) continuous. More
recently, M. Burgos and the authors of this note prove, in [13], that every
bi-orthogonality preserving linear surjection between two von Neumann algebras
(or between two compact C∗-algebras) is automatically continuous (compare
[40], [41] for recent additional generalisations).
The main goal of this section is to describe the orthogonality preserving
linear mappings between $C(K)^{\tau}$-spaces. Among the consequences, we
establish that every orthogonality preserving linear bijection between unital
commutative real C∗-algebras is automatically continuous. We shall provide an
example of an orthogonality preserving linear bijection between
$C(K)^{\tau}$-spaces which is not bi-orthogonality preserving and give a
characterisation of bi-orthogonality preserving linear maps.
We shall borrow and adapt some of the ideas developed in those previously
mentioned papers (cf. [8, 9] and [30]). In order to have a good balance
between completeness and conciseness, we just give some sketch of the
refinements needed in our setting. In any case, the results presented here are
independent innovations and extensions of those proved by Beckenstein, Narici,
and Todd and Jarosz for $C(K)$-spaces.
Let $T:C(K_{1})^{\tau_{1}}\to C(K_{2})^{\tau_{2}}$ be an orthogonality
preserving linear mapping. Keeping in mind the notation in the previous
section, we write $L_{i}:=\mathcal{O}_{i}\cup F_{i},$ where $\mathcal{O}_{i}$
and $F_{i}$ are the subsets of $K_{i}$ given by Lemma 2.1. The map sending
each $f$ in $C(Ki)^{\tau_{i}}$ to its restriction to $L_{i}$ is a
C∗-isomorphism (and hence a surjective linear isometry) from
$C(Ki)^{\tau_{i}}$ onto the real C∗-algebra $C_{r}(L_{i})$ of all continuous
functions $f:L_{i}\to\mathbb{C}$ taking real values on $F_{i}$. Thus, studying
orthogonality preserving linear maps between $C(K)^{\tau}$ spaces is
equivalent to study orthogonality preserving linear mappings between the
corresponding $C_{r}(L)$-spaces.
Henceforth, we consider an orthogonality preserving (not necessarily
continuous) linear map $T:C_{r}(L_{1})\to C_{r}(L_{2})$, where $L_{1}$ and
$L_{2}$ are two compact Hausdorff spaces and each $F_{i}$ is a closed subset
of $L_{i}$. Let us consider the sets
$Z_{1}=\\{s\in L_{2}:{\delta}_{s}T\hbox{ is a non-zero bounded real-linear
mapping}\\},$ $Z_{3}=\\{s\in L_{2}:{\delta}_{s}T=0\\},\hbox{ and
}Z_{2}=L_{2}\backslash(Z_{1}\cup Z_{3}).$
It is easy to see that $Z_{3}$ is closed. Following a very usual technique
(see, for example, [8, 9, 30, 16] and [17]), we can define a continuous
support map $\varphi:Z_{1}\cup Z_{2}\to L_{1}$. More concretely, for each
$s\in Z_{1}\cup Z_{2}$, we write $\hbox{supp}(\delta_{s}T)$ for the set of all
$t\in L_{1}$ such that for each open set $U\subseteq L_{1}$ with $t\in U$
there exists $f\in C_{r}(L_{1})$ with $\hbox{coz}(f)\subseteq U$ and
$\delta_{s}(T(f))\neq 0$. Actually, following a standard argument, it can be
shown that, for each $s\in Z_{1}\cup Z_{2},$ $\hbox{supp}(\delta_{s}T)$ is
non-empty and reduces exactly to one point $\varphi(s)\in L_{1}$, and the
assignment $s\mapsto\varphi(s)$ defines a continuous map from $Z_{1}\cup
Z_{2}$ to $L_{1}$. Furthermore, the value of $T(f)$ at every $s\in Z_{1}$
depends strictly on the value $f(\varphi(s))$. More precisely, for each $s\in
Z_{1}$ with $\varphi(s)\notin F_{1}$, the value $T(g)(s)$ is the same for
every function $g\in C_{r}(L_{1})$ with $g\equiv i$ on a neighborhood of
$\varphi(s).$ Thus, defining $T(i)(s):=0$ for every $s\in Z_{3}\cup Z_{2}$ and
for every $s\in Z_{1}$ with $\varphi(s)\in F_{1}$, and $T(i)(s):=T(g)(s)$ for
every $s\in Z_{1}\cup Z_{2}$ with $\varphi(s)\notin F_{1}$, where $g$ is any
element in $C_{r}(L_{1})$ with $g\equiv i$ on a neighborhood of $\varphi(s),$
we get a (well-defined) mapping $T(i):L_{2}\to\mathbb{C}.$ It should be
noticed that “$T(i)$” is just a symbol to denoted the above mapping and not an
element in the image of $T$. In this setting, the identity
$T(f)(s)=T(1)(s)\ \Re\hbox{e}f(\varphi(s))+T(i)(s)\ \Im\hbox{m}f(\varphi(s)),$
holds for every $s\in Z_{1}.$ Clearly, $T(1)(s),T(i)(s)\in\mathbb{R}$, for
every $s\in F_{2}$ and $\left|T(1)(s)\right|+\left|T(i)(s)\right|\neq 0$, for
every $s\in Z_{1}$.
The following property also follows from the definition of $\varphi$ by
standard arguments: Under the above conditions, let $s$ be an element in
$Z_{1}\cup Z_{2},$ then
(10) $\delta_{s}T(f)=0\hbox{ for every }f\in C_{r}(L_{1})\hbox{ with
}\varphi(s)\notin\overline{\hbox{coz}(f)}.$
###### Lemma 3.1.
The mapping $T(i)$ is bounded on the set $\varphi^{-1}(\mathcal{O}_{1})$.
Furthermore, the inequality
$\left|T(f)(s)\right|\leq\|T(1)\|+\sup_{\widetilde{s}\in\varphi^{-1}(\mathcal{O}_{1})}|T(i)(\widetilde{s})|$
holds for all $s\in Z_{1}$ and all $f\in C_{r}(L_{1})$ with
$|\Re\hbox{e}(f)|,|\Im\hbox{m}(f)|\leq 1.$
###### Proof.
Arguing by contradiction, we suppose that, for each natural $n$, there exists
$s_{n}\in\varphi^{-1}(\mathcal{O}_{1})$ such that
$\left|T(i)(s_{n})\right|>n^{3}$. The elements $s_{n}^{\prime}s$ can be chosen
so that $\varphi(s_{n})\neq\varphi(s_{m})$ for $n\neq m$, and consequently we
can find a sequence of pairwise disjoint open subsets $(U_{n})$ of
$\mathcal{O}_{1}$ with $\varphi(s_{n})\in U_{n}$. It is easily seen that we
can define a function $\displaystyle g=\sum_{n=1}^{\infty}i\ g_{n}\in
C_{r}(L_{1})$ with coz$(g_{n})\subset U_{n}$, $0\leq
g_{n}\leq\frac{1}{n^{2}}$, and $g_{n}\equiv\frac{1}{n^{2}}$ on a neighborhood
of $s_{n}$, for all $n$. By the form of $g$, and since $T$ is orthogonality
preserving, we have $|T(g)(s_{n})|=n^{2}|T(i)(s_{n})|>n$ for all $n$, which is
absurd. ∎
We can easily show now that $Z_{2}$ is an open subset of $L_{2}$. With this
aim, we consider an element $s_{0}$ in $Z_{2}$. We can pick a function $f\in
C_{r}(L_{1})$ such that $\|f\|\leq 1$ and
$|T(f)(s_{0})|>1+\|T(1)\|+\sup_{\widetilde{s}\in\varphi^{-1}(\mathcal{O}_{1})}|T(i)(\widetilde{s})|.$
Since
$\displaystyle\left|T(f)(s)\right|\leq\|T(1)\|+\sup_{\widetilde{s}\in\varphi^{-1}(\mathcal{O}_{1})}|T(i)(\widetilde{s})|<|T(f)(s_{0})|-1$,
for every $s\in Z_{1}\cup Z_{3},$ we conclude that there exists an open
neighborhood of $s_{0}$ contained in $Z_{2}$.
The next theorem resumes the above discussion.
###### Theorem 3.2.
In the notation above, let $T:C_{r}(L_{1})\to C_{r}(L_{2})$ be an
orthogonality preserving linear mapping. Then $L_{2}$ decomposes as the union
of three mutually disjoint subsets $Z_{1},Z_{2},$ and $Z_{3}$, where $Z_{2}$
is open and $Z_{3}$ is closed, there exist a continuous support map
$\varphi:Z_{1}\cup Z_{2}\to L_{1}$, and a bounded mapping
$T(i):L_{2}\to\mathbb{C}$ which is continuous on
$\varphi^{-1}(\mathcal{O}_{1})$ satisfying:
$T(i)(s)\in\mathbb{R}\ (\forall s\in F_{2}),\ T(i)(s)=0,\ (\forall s\in
Z_{3}\cup Z_{2}\hbox{ and }\forall s\in Z_{1}\hbox{ with }\varphi(s)\in
F_{1}),$ (11) $\left|T(1)(s)\right|+\left|T(i)(s)\right|\neq 0,\ (\forall s\in
Z_{1}),$ (12) $T(f)(s)=T(1)(s)\ \Re\hbox{e}f(\varphi(s))+T(i)(s)\
\Im\hbox{m}f(\varphi(s)),\hbox{ {\rm(}$\forall s\in Z_{1},f\in
C_{r}(L_{1}){\rm)},$}$ $T(f)(s)=0,\hbox{ {\rm(}$\forall s\in Z_{3},f\in
C_{r}(L_{1})${\rm),}}$
and for each $s\in L_{2}$, the mapping $C_{r}(L_{1})\to\mathbb{C}$, $f\mapsto
T(f(s)),$ is unbounded if, and only if, $s\in Z_{2}$. Furthermore, the set
$\varphi(Z_{2})$ is finite.
###### Proof.
Everything has been substantiated except perhaps the statement concerning the
set $\varphi(Z_{2})$. Arguing by contradiction, we assume the existence of a
sequence $(s_{n})$ in $Z_{2}$ such that $\varphi(s_{n})\neq\varphi(s_{m})$ for
every $n\neq m$. Find a sequence $(U_{n})$ of mutually disjoint open subsets
of $L_{1}$ satisfying $\varphi(s_{n})\in U_{n}$ and a sequence
$(f_{n})\subseteq C_{r}(L_{1})$ such that $\|f_{n}\|\leq\frac{1}{n}$,
$\hbox{coz}(f_{n})\subseteq U_{n}$ and $|\delta_{s_{n}}T(f_{n})|>n$, for every
$n\in\mathbb{N}$. The element $\displaystyle f=\sum_{n=1}^{\infty}f_{n}$ lies
in $C_{r}(L_{1}),$ and for each natural $n_{0}$,
$\displaystyle{f_{n_{0}}\perp\sum_{n=1,n\neq n_{0}}^{\infty}f_{n}}.$ Thus,
$|\delta_{s_{n_{0}}}T(f)|\geq|\delta_{s_{n_{0}}}T(f_{n_{0}})|>n_{0},$ which is
impossible. ∎
###### Remark 3.3.
The mapping $T(i):L_{2}\to\mathbb{C}$ has been defined to satisfy $T(i)(s)=0,$
for all $s\in Z_{3}\cup Z_{2}$ and for all $s\in Z_{1}$ with $\varphi(s)\in
F_{1}.$ It should be noticed here that the value $T(i)(s)$ is uniquely
determined only when $s\in Z_{1}$ and $\varphi(s)\notin F_{1}$. There are some
other choices for the values of $T(i)(s)$ at $s\in Z_{3}\cup Z_{2}$ and at
$s\in Z_{1}$ with $\varphi(s)\in F_{1}$ under which conditions $(\ref{eq 1 thm
32})$ and $(\ref{eq 2 thm 32})$ are satisfied.
###### Remark 3.4.
We shall now explore some of the consequences derived from Theorem 3.2. Let
$T:C_{r}(L_{1})\to C_{r}(L_{2})$ be an orthogonality preserving linear
mapping.
1. $(a)$
The set $Z_{3}$ is empty whenever $T$ is surjective.
2. $(b)$
$Z_{3}=\emptyset$ implies that $Z_{1}=L_{2}\backslash Z_{2}$ is a compact
subset of $L_{2}.$
3. $(c)$
$\varphi(Z_{2})$ is a finite set of non-isolated points in $L_{1}$. Indeed, if
$\varphi(s_{0})=t_{0}$ is isolated for some $s_{0}\in Z_{2}$, then we can find
an open set $U\subseteq L_{1}$ such that $U\cap K_{1}=\\{t_{0}\\}$. Therefore,
for each $f\in C_{r}(L_{1})$ with $f(t_{0})=0$ we have $\delta_{s_{0}}T(f)=0$.
Pick an arbitrary $h\in C_{r}(L_{1})$. Clearly, $\chi_{{}_{t_{0}}}\in
C_{r}(L_{1})$, while $i\chi_{{}_{t_{0}}}$ lies in $C_{r}(L_{1})$ if, and only
if, $t_{0}\notin F_{1}$. Therefore,
$h_{0}=\Re\hbox{e}(h(t_{0}))\chi_{{}_{t_{0}}}+\Im\hbox{m}(h(t_{0}))\
i\chi_{{}_{t_{0}}}$
lies in $C_{r}(L_{1})$ and $(h-h_{0})(t_{0})=0$.
Assume first that $t_{0}\notin F_{1}$. Denoting
$\lambda_{0}=\delta_{s_{0}}T(\chi_{{}_{t_{0}}})$ and
$\mu_{{}_{0}}=\delta_{s_{0}}T(i\chi_{{}_{t_{0}}}),$ we have
$\delta_{s_{0}}T(h)=\delta_{s_{0}}T(h_{0})=\lambda_{0}\Re\hbox{e}(h(t_{0}))+\mu_{{}_{0}}\Im\hbox{m}(h(t_{0}))$
$=\frac{\lambda_{0}-i\mu_{{}_{0}}}{2}\delta_{t_{0}}(h)+\frac{\lambda_{0}+i\mu_{{}_{0}}}{2}\overline{\delta_{t_{0}}}(h).$
This shows that
$\delta_{s_{0}}T=\frac{\lambda_{0}-i\mu_{{}_{0}}}{2}\delta_{t_{0}}+\frac{\lambda_{0}+i\mu_{{}_{0}}}{2}\overline{\delta_{t_{0}}}$
is a continuous mapping from $C_{r}(L_{1})$ to $\mathbb{C}$, which is
impossible.
When $t_{0}\in F_{1}$ we have $\delta_{s_{0}}T=\lambda_{0}\delta_{t_{0}}$ is a
continuous mapping from $C_{r}(L_{1})$ to $\mathbb{R}$, which is also
impossible.
4. $(d)$
$T$ surjective implies
$\varphi(Z_{1}\cap\mathcal{O}_{2})\subseteq\mathcal{O}_{1}$. Suppose, on the
contrary that there exists $s_{0}\in Z_{1}\cap\mathcal{O}_{2}$ with
$\varphi(s_{0})\in F_{1}$. By $(\ref{eq 2 thm 32})$,
$T(f)(s_{0})=T(1)(s_{0})\Re\hbox{e}f(\varphi(s_{0})),$
for every $f\in C_{r}(L_{1})$. It follows from the surjectivity of $T$,
together with the condition $s_{0}\in\mathcal{O}_{2},$ that for every complex
number $\omega$ there exists a real $\lambda$ satisfying
$\omega=T(1)(s_{0})\lambda,$ which is impossible.
5. $(e)$
Suppose $T$ is surjective and fix $s_{0}\in Z_{1}\cap\mathcal{O}_{2}$. The
mapping $\delta_{s_{0}}T$ is a bounded real-linear mapping from $C_{r}(L_{1})$
onto $\mathbb{C}$. On the other hand, by $(\ref{eq 2 thm 32})$,
$\delta_{s_{0}}T(f)=T(1)(s_{0})\Re\hbox{e}f(\varphi(s_{0}))+T(i)(s_{0})\Im\hbox{m}f(\varphi(s_{0})),\
(\forall f\in C_{r}(L_{1})).$
Thus, $T$ being surjective implies that the space
$\mathbb{C}_{\mathbb{R}}=\mathbb{R}\times\mathbb{R}$ is linearly spanned by
the elements $T(1)(s_{0})$ and $T(i)(s_{0})$. Therefore, for each $s_{0}\in
Z_{1}\cap\mathcal{O}_{2},$ the set $\\{T(1)(s_{0}),T(i)(s_{0})\\}$ is a basis
of $\mathbb{C}_{\mathbb{R}}=\mathbb{R}\times\mathbb{R}.$ Consequently, when
$T$ is surjective and $s_{0}\in Z_{1}\cap\mathcal{O}_{2},$ the condition
$T(f)(s_{0})=0$ implies $f(\varphi(s_{0}))=0.$ For any other $s_{1}\in
Z_{1}\cap\mathcal{O}_{2}$ with $\varphi(s_{0})=\varphi(s_{1})$, we have:
$T(f)(s_{0})=0\Rightarrow f(\varphi(s_{0}))=0\Rightarrow T(f)(s_{1})=0.$
The fact that $C_{r}(L_{2})$ separates points implies that $s_{1}=s_{0}$.
Thus, $\varphi$ is injective on $Z_{1}\cap\mathcal{O}_{2}$.
We can now state the main result of this section which affirms that every
orthogonality preserving linear bijection between unital commutative real
C∗-algebras is (automatically) continuous.
###### Theorem 3.5.
Every orthogonality preserving linear bijection between unital commutative
(real) C∗-algebras is (automatically) continuous.
###### Proof.
Since $T$ is surjective, $Z_{3}=\emptyset$, and hence $Z_{1}=L_{2}\backslash
Z_{2}$ is a compact subset of $L_{2}.$ It is also clear that $\varphi(L_{2})$
is compact. We claim that $\varphi(L_{2})=L_{1}.$ Otherwise, there would exist
a non-zero function $f\in C_{r}(L_{1})$ with
$\overline{\hbox{coz}(f)}\subseteq L_{1}\backslash\varphi(L_{2})$. Thus, by
$(\ref{eq varphi notin cozero})$, $T(f)=0,$ contradicting the injectivity of
$T$. By Remark 3.4$(c)$,
$\varphi(Z_{1})=\overline{\varphi(Z_{1})}=\varphi(L_{2})=\varphi(Z_{1})\cup\varphi(Z_{2})=L_{1}.$
We next see that $Z_{2}=\emptyset.$ Otherwise we can take $g\in C_{r}(L_{2})$
with $\emptyset\neq coz(g)\subset Z_{2}.$ Let $h=T^{-1}(g).$ Obviously
$Th(s)=0$ whenever $s\in Z_{1}.$ We claim that $h(t)=0,$ for every
$t\in\varphi(Z_{1})\setminus\varphi(Z_{2}).$ Let us fix
$t\in\varphi(Z_{1})\setminus\varphi(Z_{2}).$ Since $\varphi(Z_{2})$ is a
finite set there are disjoint open sets $U_{1},U_{2}$ such that $t\in
U_{1},\varphi(Z_{2})\subset U_{2}.$ Let $f\in C(L_{1},{\mathbb{R}})$ be such
that $f(t)\neq 0$ and $\overline{coz(f)}\subset U_{1}.$ We see that $T(fh)=0.$
Indeed, let $s\in L_{2}=Z_{1}\cup Z_{2}.$ If $s$ lies in $Z_{1},$ then the
maps $fh$ and $f(\varphi(s))h$ lie in $C_{r}(L_{1})$ and coincide at
$\varphi(s).$ Since $T$ is linear over ${\mathbb{R}}$ and $f$ takes real
values, we deduce, by $(\ref{eq 2 thm 32})$, that
$T(fh)(s)=f(\varphi(s))Th(s)=0.$ If $s\in Z_{2}$ then, since
$\varphi(s)\notin\overline{coz(fh)}$, then $\delta_{s}T(fh)=T(fh)(s)=0$.
We have shown that $T(fh)=0$. Thus, since $T$ is injective, $fh=0$ and
therefore $h(t)=0.$ We have therefore proved that
$coz(h)\subset\varphi(Z_{2})$ which is a finite set. This means that $h$ must
be a finite linear combination of characteristic function on points of
$\varphi(Z_{2})$ and these points must be isolated which is impossible, since
by $c)$ in Remark 3.4 no point in $\varphi(Z_{2})$ can be isolated. We have
proved that $Z_{2}=\emptyset.$ Now the fact that $T$ is continuous follows
easily. ∎
The above theorem is the first step toward extending, to the real setting,
those results proved in [30], [6], [13], [40], [38] and [48] for (complex)
C∗-algebras.
Orthogonality preserving linear bijections enjoy an interesting additional
property.
###### Proposition 3.6.
In the notation of this section, let $T:C_{r}(L_{1})\to C_{r}(L_{2})$ be an
orthogonality preserving linear bijection. Then $T^{-1}$ preserves invertible
elements, that is, $T^{-1}(g)$ is invertible whenever $g$ is an invertible
element in $C_{r}(L_{2})$.
###### Proof.
Take an invertible element $g\in C_{r}(L_{2})$. Let $f$ be the unique element
in $C_{r}(L_{1})$ satisfying $T(f)=g$. Theorem 3.2 implies that
$0\neq g(s)=T(f)(s)=T(1)(s)\ \Re\hbox{e}f({\varphi(s)})+T(i)(s)\
\Im\hbox{m}f({\varphi(s)}),$
for every $s\in Z_{1}$. This assures that $f({\varphi(s)})\neq 0,$ for every
$s\in Z_{1},$ and since $\varphi(Z_{1})=L_{1},$ $f=T^{-1}(g)$ must be
invertible in $C_{r}(L_{1})$. ∎
In the setting of complex Banach algebras, it follows from the Gleason-Kahane-
Żelazko theorem that a linear transformation $\phi$ from a unital,
commutative, complex Banach algebra $A$ into $\mathbb{C}$ satisfying
$\phi(1)=1$ and $\phi(a)\neq 0$ for every invertible element $a$ in $A$ is
multiplicative, that is, $\phi(ab)=\phi(a)\phi(b)$ (see [20, 32]). Although,
the Gleason-Kahane-Żelazko theorem fails for real Banach algebras, S.H.
Kulkarni found in [35] the following reformulation: a linear map $\phi$ from a
real unital Banach algebra ${A}$ into the complex numbers is multiplicative if
$\varphi(1)=1$ and $\phi(a)^{2}+\phi(b)^{2}\neq 0$ for every $a,b\in{A}$ with
$ab=ba$ and $a^{2}+b^{2}$ invertible. It is not clear that statement $(b)$ in
the above proposition can be improved to get the hypothesis of Kulkarni’s
theorem. The structure of orthogonality preserving linear mappings between
$C_{r}(L)$-spaces described in Theorem 3.2 invites us to affirm that they are
not necessarily multiplicative.
### 3.1. Bi-orthogonality preservers
As a consequence of the description of orthogonality preserving linear maps
given in [30], it can be shown that an orthogonality preserving linear
bijection between (complex) $C(K)$-spaces is bi-orthogonality preserving. It
is natural to ask wether every orthogonality preserving linear bijection
between commutative (unital) real C∗-algebras is bi-orthogonality preserving.
This is known to be true in two cases: first, between spaces
$C_{\mathbb{R}}(K)$ of real (and also complex) valued functions on a compact
Hausdorff space $K$, as it is well-known; second, between spaces of the type
$C_{\mathbb{R}}(K;\mathbb{R}^{n})$ (compare [16, Section 3]). Spaces like
those we are dealing with in this paper need not satisfy this property, that
is, there exists an orthogonality preserving linear bijection
$T:C_{r}(L_{1})\to C_{r}(L_{2})$ which is not bi-orthogonality preserving (and
even $L_{1}$ and $L_{2}$ are not homeomorphic either).
###### Example 3.7.
Let $L_{1}=\\{t_{1},t_{2},t_{3}\\}$ $L_{2}=\\{s_{1},s_{2},s_{3},s_{4}\\}$ with
$\mathcal{O}_{1}=\\{t_{1},t_{3}\\}$, $\mathcal{O}_{2}=\\{s_{1}\\}$,
$F_{1}=\\{t_{2}\\}$ and $F_{2}=\\{s_{2},s_{3},s_{4}\\}$. Define
$\varphi:L_{2}\to L_{1}$ by $\varphi(s_{i})=t_{i},$ for $i=1,2,$ and
$\varphi(s_{i})=t_{3},$ for $i=3,4$. It is easy to check that
$T(f)(s_{i})=f(\varphi(s_{i}))$ if $i=1,2$, and
$T(f)(s_{3})=\Re\hbox{e}f(t_{3}),$ $T(f)(s_{4})=\Im\hbox{m}f(t_{3})$ is an
orthogonality preserving linear bijection, but $T^{-1}$ is not orthogonality
preserving.
In the above example, $\varphi^{-1}(\mathcal{O}_{1})\cap F_{2}$ is non-empty.
Our next result shows that a topological condition on $F_{2}$ assures that an
orthogonality preserving linear bijection between unital commutative real
C∗-algebras is bi-orthogonality preserving.
###### Proposition 3.8.
In the notation of this section, let $T:C_{r}(L_{1})\to C_{r}(L_{2})$ be an
orthogonality preserving linear bijection (not assumed to be bounded). The
following statements hold:
1. $(a)$
If $T$ is bi-orthogonality preserving then $\varphi:L_{2}\to L_{1}$ is a
(surjective) homeomorphism, $\varphi(F_{2})=F_{1},$ and
$\varphi(\mathcal{O}_{2})=\mathcal{O}_{1}$. In particular,
$\varphi^{-1}(\mathcal{O}_{1})\cap F_{2}=\emptyset$.
2. $(b)$
If $F_{2}$ has empty interior then $T$ is biorthogonality preserving.
###### Proof.
$(a)$ If $T$ is bi-orthogonality preserving, it can be easily seen that
$\varphi:L_{2}\to L_{1}$ is a homeomorphism, and for each $s\in L_{2}$,
supp$(\delta_{\varphi(s)}T^{-1})=\\{s\\}$. By Remark 3.4$(d)$, applied to $T$
and $T^{-1}$, we have $\varphi(F_{2})=F_{1}$ and
$\varphi(\mathcal{O}_{2})=\mathcal{O}_{1}$. Then
$\varphi^{-1}(\mathcal{O}_{1})\cap F_{2}=\emptyset.$ So, $a)$ is clear.
$(b).$ Let us assume that $F_{2}$ has empty interior. Arguing by contradiction
we suppose that $T^{-1}$ is not orthogonality preserving. Then there exist
$f_{1},f_{2}\in C_{r}(L_{1})$ with $f_{1}f_{2}\neq 0$, but $T(f_{1})\perp
T(f_{2})$. Thus $U:=\hbox{coz}(f_{1})\cap\hbox{coz}(f_{2})$ is a non-empty
open subset of $L_{1}.$ Keeping again the notation of Theorem 3.2 for $T$, we
recall that, by Theorem 3.5 and Remark 3.4, $Z_{3}=\emptyset,$
$Z_{2}=\emptyset$, $\varphi(L_{2})=L_{1}$,
$\varphi(\mathcal{O}_{2})\subset\mathcal{O}_{1},$ $\varphi|_{\mathcal{O}_{2}}$
is injective, and for each $s\in\mathcal{O}_{2}$, and $\\{T(1)(s),T(i)(s)\\}$
is a basis of $\mathbb{C}_{\mathbb{R}}=\mathbb{R}\times\mathbb{R}.$
By the form of $T$, there are no points of $\varphi(\mathcal{O}_{2})$ in
$U=\hbox{coz}(f_{1})\cap\hbox{coz}(f_{2})$ (because for each
$s\in\mathcal{O}_{2},$ $T(f)(s)\neq 0$ when $f(\varphi(s))\neq 0$). Now, let
$k$ be a non-zero element in $C(L_{1},\mathbb{R})$, with
coz$(k)\subseteq\hbox{coz}(f_{1})\cap\hbox{coz}(f_{2})$. By Theorem 3.2
$(\ref{eq 2 thm 32})$, it is clear that
$\varphi(\hbox{coz}(T(k)))\subseteq\hbox{coz}(k)$, and hence, since
$\varphi(\mathcal{O}_{2})\subseteq\mathcal{O}_{1}$, $\hbox{coz}(T(k))$ is a
non-empty subset of $F_{2}$, against our hypotheses. ∎
As we have already seen, an orthogonality preserving linear bijection between
$C_{r}(L)$-spaces needs not to be biorthogonality preserving. Example 3.7 also
shows that, unlike in the complex case, the existence of an orthogonality
preserving linear bijection between $C_{r}(L)$-spaces does not guarantee that
the corresponding compacts spaces are homeomorphic. We next provide a
characterisation of those (linear) mappings which are bi-orthogonality
preserving. As a consequence, we shall see that if there exists a bi-
orthogonality preserving linear map $T:C_{r}(L_{1})\to C_{r}(L_{2})$ then
$L_{1}$ and $L_{2}$ are homeomorphic.
###### Theorem 3.9.
Let $T:C_{r}(L_{1})\to C_{r}(L_{2})$ be a mapping. The following statements
are equivalent:
1. $(a)$
$T$ is a bi-orthogonality preserving linear surjection;
2. $(b)$
There exists a (surjective) homeomorphism $\varphi:L_{2}\to L_{1}$ with
$\varphi(\mathcal{O}_{2})=\mathcal{O}_{1},$ a function
$a_{1}=\gamma_{1}+i\gamma_{2}$ in $C_{r}(L_{2})$ with $a_{1}(s)\neq 0$ for all
$s\in L_{2},$ and a function $a_{2}=\eta_{1}+i\eta_{2}:L_{2}\to{\mathbb{C}}$
continuous on $\mathcal{O}_{2}$ with the property that
$0<\inf_{s\in\mathcal{O}_{2}}\left|\det\left(\begin{array}[]{cc}\gamma_{1}(s)&\eta_{1}(s)\\\
\gamma_{2}(s)&\eta_{2}(s)\\\
\end{array}\right)\right|\leq\sup_{s\in\mathcal{O}_{2}}\left|\det\left(\begin{array}[]{cc}\gamma_{1}(s)&\eta_{1}(s)\\\
\gamma_{2}(s)&\eta_{2}(s)\\\ \end{array}\right)\right|<+\infty,$
such that
$T(f)(s)=a_{1}(s)\ \Re\hbox{e}f({\varphi(s)})+a_{2}(s)\
\Im\hbox{m}f({\varphi(s)})$
for all $s\in L_{2}$ and $f\in C_{r}(L_{1}).$
###### Proof.
$(a)\Rightarrow(b)$. Since every bi-orthogonality preserving linear mapping is
injective, we can assume that $T:C_{r}(L_{1})\to C_{r}(L_{2})$ is a bi-
orthogonality preserving linear bijection. We keep the notation given in
Theorem 3.2. We have already shown that $Z_{3}=\emptyset,$ $Z_{2}=\emptyset$,
$\varphi:L_{2}\to L_{1}$ is a surjective homeomorphism,
$\varphi(\mathcal{O}_{2})=\mathcal{O}_{1},$ and for each
$s\in\mathcal{O}_{2}$, $\\{T(1)(s),T(i)(s)\\}$ is a basis of
$\mathbb{C}_{\mathbb{R}}=\mathbb{R}\times\mathbb{R}$ (compare Theorem 3.5,
Remark 3.4 and Proposition 3.8). Taking $a_{1}=T(1)=\gamma_{1}+i\gamma_{2}$
and $a_{2}=T(i)=\eta_{1}+i\eta_{2}$ we only have to show that
$0<\inf_{s\in\mathcal{O}_{2}}\left|\det\left(\begin{array}[]{cc}\gamma_{1}(s)&\eta_{1}(s)\\\
\gamma_{2}(s)&\eta_{2}(s)\\\
\end{array}\right)\right|\leq\sup_{s\in\mathcal{O}_{2}}\left|\det\left(\begin{array}[]{cc}\gamma_{1}(s)&\eta_{1}(s)\\\
\gamma_{2}(s)&\eta_{2}(s)\\\ \end{array}\right)\right|<+\infty.$
Let us denote $M_{s}=\left(\begin{array}[]{cc}\gamma_{1}(s)&\eta_{1}(s)\\\
i\gamma_{2}(s)&i\eta_{2}(s)\\\ \end{array}\right).$ Clearly $\det(M_{s})\neq
0$, for every $s\in\mathcal{O}_{2}$ and
$T(f)(s)=M_{s}\cdot\left(\begin{array}[]{c}\Re\hbox{e}f(\varphi(s))\\\
\Im\hbox{m}f(\varphi(s))\\\ \end{array}\right),$ for every $f\in
C_{r}(L_{1}),s\in L_{2}$. By the boundedness of $T(1):L_{2}\to\mathbb{C}$ and
$T(i)|_{\mathcal{O}_{2}}:\mathcal{O}_{2}\to\mathbb{C}$ (see Lemma 3.1) there
exists $M>0$ such that $\left|det(M_{s})\right|\leq M$ for all
$s\in\mathcal{O}_{2}.$
Applying the above arguments to the mapping $T^{-1}$ we find a surjective
homeomorphism $\psi=\varphi^{-1}:L_{1}\to L_{2}$, a mapping
$T^{-1}(i):L_{1}\to L_{2}$ and $m>0,$ such that
$\psi(\mathcal{O}_{1})=\mathcal{O}_{2}$, for each $t\in\mathcal{O}_{1}$,
$\\{T^{-1}(1)(t),T^{-1}(i)(t)\\}$ is a basis of
$\mathbb{C}_{\mathbb{R}}=\mathbb{R}\times\mathbb{R},$
$T^{-1}(g)(t)=N_{t}\cdot\left(\begin{array}[]{c}\Re\hbox{e}g(\psi(t))\\\
\Im\hbox{m}g(\psi(t))\\\ \end{array}\right)$ $(g\in C_{r}(L_{2}),t\in L_{1})$,
$\left|\det(N_{t})\right|\leq m,$ for all $t\in\mathcal{O}_{1},$ where
$N_{t}=\left(\begin{array}[]{cc}\Re\hbox{e}T^{-1}(1)(t)&\Re\hbox{e}T^{-1}(i)(t)\\\
i\Im\hbox{m}T^{-1}(1)(t)&i\Im\hbox{m}T^{-1}(i)(t)\\\ \end{array}\right).$ It
can be easily seen that, for each $s\in\mathcal{O}_{2},$
$N_{\varphi(s)}=M_{s}^{-1},$ which shows that
$\left|det(M_{s})\right|\geq\frac{1}{m}$, for all $s\in\mathcal{O}_{2}.$
$(b)\Rightarrow(a)$. Let $T:C_{r}(L_{1}):\to C_{r}(L_{2})$ be a mapping
satisfying the hypothesis in $(b).$ Clearly, $T$ is linear, and since
$\varphi(F_{2})=F_{1},$ $Tf(s)\in{\mathbb{R}}$ for all $s\in F_{2}$ and $f\in
C_{r}(L_{1})$ (that is, $T(f)\in C_{r}(L_{2})$). We can easily check that,
under these hypothesis, $T$ is injective and preserves orthogonality.
We shall now prove that $T$ is surjective. Indeed, for each
$s\in\mathcal{O}_{2}$
$T(f)(s)=\left(\begin{array}[]{c}\Re\hbox{e}g(s)\\\ \Im\hbox{m}g(s)\\\
\end{array}\right)=\left(\begin{array}[]{cc}\gamma_{1}(s)&\eta_{1}(s)\\\
i\gamma_{2}(s)&i\eta_{2}(s)\\\
\end{array}\right)\cdot\left(\begin{array}[]{c}\Re\hbox{e}f({\varphi(s)})\\\
\Im\hbox{m}f({\varphi(s)})\\\ \end{array}\right)$
$=M_{s}\cdot\left(\begin{array}[]{c}\Re\hbox{e}f({\varphi(s)})\\\
\Im\hbox{m}f({\varphi(s)})\\\ \end{array}\right),$
thus,
$\left(\begin{array}[]{c}\Re\hbox{e}f({\varphi(s)})\\\
\Im\hbox{m}f({\varphi(s)})\\\
\end{array}\right)=M_{s}^{-1}\cdot\left(\begin{array}[]{c}\Re\hbox{e}g(s)\\\
\Im\hbox{m}g(s)\\\ \end{array}\right).$
We define $b_{1}(t):L_{1}\to{\mathbb{C}}$ and
$b_{2}:\mathcal{O}_{1}\to{\mathbb{C}}$ by
$b_{1}(t)=\widetilde{\gamma}_{1}(t)+i\widetilde{\gamma}_{2}(t)$ and
$b_{2}=\widetilde{\eta}_{1}(t)+i\widetilde{\eta}_{2}(t)$
($t\in\mathcal{O}_{1}$), where
$M_{\varphi^{-1}(t)}^{-1}=\left(\begin{array}[]{cc}\widetilde{\gamma}_{1}(t)&\widetilde{\eta}_{1}(t)\\\
i\widetilde{\gamma}_{2}(t)&i\widetilde{\eta}_{2}(t)\\\ \end{array}\right),$
and $b_{1}(t)=\frac{1}{\gamma_{1}(\varphi^{-1}(t))},$ for every $t\in F_{1}$.
Then $S:C_{r}(L_{2})\to C_{r}(L_{1}),$ defined by $S(g)(t)=b_{1}(t)\
\Re\hbox{e}g(\varphi^{-1}(t))+b_{2}(t)\ \Im\hbox{m}g(\varphi^{-1}(t)),$ is
linear, preserves orthogonality and it is easy to check that $S=T^{-1}.$ It
follows that $T$ is bi-orthogonality preserving. ∎
Let $T$ be a bi-orthogonality preserving linear mapping with associated
homeomorphism $\varphi:L_{2}\to L_{1}$. Clearly, the operator
$S:C_{r}(L_{1})\to C_{r}(L_{2})$, $S(f)(s):=f(\varphi(s))$ is a ∗-isomorphism.
Having in mind that a linear mapping $T:A\to B$ between real C∗-algebras is a
∗-isomorphism if, and only if, the complex linear extension
$\widetilde{T}:A\oplus iA\to B\oplus iB$, $\widetilde{T}(a+ib)=T(a)+iT(b)$ is
a ∗-isomorphism, we get the following corollary.
###### Corollary 3.10.
The following statements are equivalent:
1. $(a)$
There exists a bi-orthogonality preserving linear bijection $T:C_{r}(L_{1})\to
C_{r}(L_{2})$;
2. $(b)$
There exists a C∗-isomorphism $S:C_{r}(L_{1})\to C_{r}(L_{2})$;
3. $(c)$
There exists a C∗-isomorphism $\widetilde{S}:C(L_{1})\to C(L_{2})$;
4. $(d)$
$L_{1}$ and $L_{2}$ are homeomorphic.$\hfill\Box$
Acknowledgements: The authors gratefully thank to the Referee for the
constructive comments and detailed recommendations which definitely helped to
improve the readability and quality of the paper.
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|
arxiv-papers
| 2013-09-16T07:32:31 |
2024-09-04T02:49:50.953116
|
{
"license": "Public Domain",
"authors": "Jorge J. Garc\\'es and Antonio M. Peralta",
"submitter": "Antonio M. Peralta",
"url": "https://arxiv.org/abs/1309.3839"
}
|
1309.4074
|
Three-Dimensional Smoothed Particle Hydrodynamics
Method for Simulating Free Surface Flows
Rizal Dwi Prayogoa,b, Christian Fredy Naaa
aFaculty of Mathematics and Natural Sciences, Institut Teknologi Bandung, Jl.
Ganesha 10, Bandung 40132 Indonesia, E-mail: [email protected],
[email protected]
bGraduate School of Natural Science and Technology, Kanazawa University,
Kakuma, Kanazawa 920-1192 Japan,
Abstract. In this paper, we applied an improved Smoothing Particle
Hydrodynamics (SPH) method by using gradient kernel renormalization in three-
dimensional cases. The purpose of gradient kernel renormalization is to
improve the accuracy of numerical simulation by improving gradient kernel
approximation. This method is implemented for simulating free surface flows,
in particular dam break case with rigid ball structures and the propagation of
waves towards a slope in a rectangular tank.
Keywords: Smoothed particle hydrodynamics, free surface flows, gradient kernel
renormalization
## 1 Introduction
Computational Fluid Dynamics (CFD) using Smoothed Particle Hydrodynamics (SPH)
has a wide range of applications to solve problem in engineering and science.
SPH is a mesh-free Lagrangian method and well suited to the simulation of
complex and free surface flows. The SPH method was originally used to model
astrophysical problems by Lucy [7] and Gingold and Monaghan [8]. Three-
dimensional SPH method has been studied numerically by Monaghan [6] in a field
of astrophysical fluid dynamics processes.
We obtain the SPH equations from the continuum equations of fluid dynamics by
interpolating from set of points which may be disordered. This interpolation
is based on the theory of integral interpolants using interpolation kernels
which approximate the delta function. The interpolants being analytic
functions can be differentiated without using grids.
Monaghan [4] studied the application of the particle method SPH to free
surface problems in two-dimensional cases. In this paper, we applied an
improved SPH method by using gradient kernel renormalization in three-
dimensional cases. The purpose of gradient kernel renormalization is to
improve the accuracy of the simulations [1].
In the following, first the general concept of SPH method is given. The
improved SPH method using gradient kernel renormalization is introduced and
described in detail. This improved SPH method is implemented for simulating
free surface flows, in particular dam break case with rigid ball structures
and the propagation of waves towards a slope in a rectangular tank.
## 2 Smoothed particle hydrodynamics
The SPH equations are described in detail by Liu and Liu [2]. In this paper,
we consider the application of three-dimensional SPH to free surface problems.
The SPH method represents continuous fluid using a set of particles. Each
particle $i$ has physical quantities, such as mass $m_{i}$, position
$\mathbf{r}_{i}$, velocity $\mathbf{v}_{i}$, density $\rho_{i}$, and pressure
$P_{i}$. Each particle in the SPH method is associated with a support domain.
The SPH approximation, which consists of the particle approximation and the
kernel approximation, is performed within the current support domain. The
value of a function defining a physical quantity can be approximated by its
values at a number of neighboring particles. The SPH method uses the concept
of integral representation of a field function $f(x)$ by the following
identity
$<f(\mathbf{r})>=\int_{\Omega}f(\mathbf{r^{\prime}})W(|\mathbf{r}-\mathbf{r^{\prime}}|,h)d\mathbf{r^{\prime}},$
(1)
where $\mathbf{r}$ and $\mathbf{r}^{\prime}$ are the position vectors, $W$ is
the smoothing function or kernel function, and $h$ is the smoothing length
defining the influence radius of $W$.
Figure 1: The influence radius of $W$.
In SPH approximation, there are various kernel functions. Since it is affects
the accuracy and stability of numerical results, the choice of kernel function
$W$ is important to consider. The integral representation should satisfies
several conditons. That is the normalization condition
$\int_{\Omega}W(|\mathbf{r}-\mathbf{r^{\prime}}|,h)d\Omega=1,$
delta function property
$\lim_{h\to
0}W(|\mathbf{r}-\mathbf{r^{\prime}}|,h)=\delta(|\mathbf{r}-\mathbf{r^{\prime}}|),$
moreover, often the compact support condition is required
$W(|\mathbf{r}-\mathbf{r^{\prime}}|,h)=0\text{ outside of support domain.}$
We converted the continuous integral representation (1) into discretized forms
as a summation over all the particles in the support domain. This process is
also commonly known as particle approximation in the SPH literature [2].
Writing the particle approximation as follows
$<f(\mathbf{r})>=\sum_{j=1}^{N}\dfrac{m_{j}}{\rho_{j}}f(\mathbf{r_{j}})W(|\mathbf{r}-\mathbf{r_{j}}|,h),$
where $m_{j}$ and $\rho_{j}$ are the mass and density of the particle $j$,
respectively, and $j=1,2,\ldots,N$, where $N$ is the total number of
neighboring particles in the influence domain $\Omega$. In this paper, we use
cubic spline kernel as follows [5]
$W(q,h)=\dfrac{1}{\pi
h^{3}}\begin{cases}1-\frac{3}{2}q^{2}+\frac{3}{4}q^{3},&\text{$0\leq q<1$}\\\
\frac{1}{4}(2-q)^{3},&\text{$1\leq q<2$}\\\ 0,&\text{otherwise,}\end{cases}$
where $q=\dfrac{|\mathbf{r_{i}}-\mathbf{r_{j}}|}{h}$ is the relative distance
of particle $i$ and $j$.
## 3 Numerical model
### 3.1 The continuity equation
The continuity equation is based on the conservation of mass. We write the
continuity equation in the form
$\dfrac{D\rho}{Dt}=-\rho\nabla\mathbf{v},$ (2)
where $\mathbf{v}$ and $\rho$ are velocity and density, respectively. Writing
(2) in SPH discretization form as in [3], we obtain
$\dfrac{D\rho_{i}}{Dt}=\rho_{i}\sum_{j}\dfrac{m_{j}}{\rho_{j}}(\mathbf{v_{i}}-\mathbf{v_{j}})\nabla_{r_{i}}W(|\mathbf{r_{i}}-\mathbf{r_{j}}|,h),$
(3)
where $\rho_{k}$ and $\mathbf{v}_{k}$ are density and velocity of particle $k$
(evaluated at $k=i$ or $k=j$), respectively, $m_{j}$ is mass of particle $j$
and
$\nabla_{r_{i}}W(|\mathbf{r_{i}}-\mathbf{r_{j}}|,h)=\dfrac{\mathbf{r_{i}}-\mathbf{r_{j}}}{|\mathbf{r_{i}}-\mathbf{r_{j}}|}\dfrac{\partial
W}{\partial r}.$
### 3.2 The momentum equation
The momentum equation is based on the conservation of momentum which is given
by
$\dfrac{D\mathbf{v}}{Dt}=-\dfrac{1}{\rho}\nabla P+\mathbf{F},$ (4)
where $\mathbf{v},\rho$ and $P$ are velocity, density, and pressure,
respectively. Here, $F$ is external force, in this case gravitational
acceleration. Writing (4) in SPH discretization form as in [3], we get
$\dfrac{D\mathbf{v_{i}}}{Dt}=-\sum_{j}m_{j}\Big{(}\dfrac{P_{i}+P_{j}}{\rho_{i}\rho_{j}}+\Pi_{ij}\Big{)}\nabla_{r_{i}}W(|\mathbf{r_{i}}-\mathbf{r_{j}}|,h)+\mathbf{F},$
(5)
where $P_{k}$ is pressure of particle $k$ (evaluated at $k=i$ or $k=j$).
In SPH, there are various formulations for viscosity. In the momentum
equation, the introduction of a viscous term is necessary not only to consider
viscid fluids and no slip boundary conditions, but also to provide the
stability to the system and to prevent inter-particle penetration. The
artificial viscosity term $\Pi_{ij}$ is added to pressure terms within the
momentum equation (5). The artificial viscosity $\Pi_{ij}$ has the form [6]
$\Pi_{ij}=\begin{cases}\dfrac{-\alpha
c\mu_{ij}+\beta\mu_{ij}^{2}}{(\rho_{i}+\rho_{j})/2},&\text{$(\mathbf{v_{i}-\mathbf{v_{j}}}).(\mathbf{r_{i}-\mathbf{r_{j}}})<0$}\\\
0,&\text{$(\mathbf{v_{i}-\mathbf{v_{j}}}).(\mathbf{r_{i}-\mathbf{r_{j}}})>0$}\end{cases}$
where
$\displaystyle\mu_{ij}=\dfrac{h(\mathbf{r_{i}-\mathbf{r_{j}}}).(\mathbf{v_{i}-\mathbf{v_{j}}})}{|\mathbf{r_{i}}-\mathbf{r_{j}}|^{2}+\eta^{2}}.$
In these expressions, $c$ is the speed of sound, $\eta=0.001$, $\alpha$ and
$\beta$ represent shear and bulk viscosity, respectively. For the problems
described here, we choose $\alpha=0.03$ and $\beta=0$.
### 3.3 The equation of state
The equation of state is used to relate density to pressure. In this paper,
the Tait’s equation of state has the form
$p=\dfrac{\rho_{0}c^{2}}{\gamma}\Big{[}\Big{(}\dfrac{\rho}{\rho_{0}}\Big{)}^{\gamma}-1\Big{]},$
where $c,\rho_{0}$, and $\gamma$ are the speed of sound, density reference,
and the polytropic constant, respectively. Note that $\gamma=7$ is usually
used for water simulations. The speed of sound $c$ is approximately
$\sqrt{100gH}$ and it is chosen in respect of a low Mach number ($Ma<0.1$) to
ensure low compressibility effects [6].
## 4 Improvement of the SPH method
In this paper, we applied an improvement to the standard SPH method by using
renormalization. This technique is to improve the accuracy of the method [1].
### 4.1 Gradient kernel renormalization
The velocity gradient in (3) can be approached by using
$\displaystyle\nabla\mathbf{v}=\nabla\mathbf{v}-\mathbf{v}\nabla 1.$
We can generalize this approach for any field $f$ by using
$\nabla f=\nabla f-f\nabla 1$ (6)
and transforming (6) into its continuous convoluted form we have
$\langle\nabla f(\mathbf{r})\rangle=\int_{\Omega}f(\mathbf{r^{\prime}})\nabla
Wd\mathbf{r^{\prime}}-f(\mathbf{r})\int_{\Omega}\nabla Wd\mathbf{r^{\prime}}.$
We recall the second order Taylor expansion
$\displaystyle\int_{\Omega}f(\mathbf{r^{\prime}})\nabla Wd\mathbf{r^{\prime}}$
$\displaystyle=f(\mathbf{r})\int_{\Omega}\nabla
Wd\mathbf{r^{\prime}}+\dfrac{\partial
f(\mathbf{r})}{\partial\mathbf{r_{1}}}\underbrace{\int_{\Omega}(\mathbf{r_{1}^{\prime}}-\mathbf{r_{1}})\nabla
Wd\mathbf{r^{\prime}}}_{A}+$ $\displaystyle\dfrac{\partial
f(\mathbf{r})}{\partial\mathbf{r_{2}}}\underbrace{\int_{\Omega}(\mathbf{r_{2}^{\prime}}-\mathbf{r_{2}})\nabla
Wd\mathbf{r^{\prime}}}_{B}+\dfrac{\partial
f(\mathbf{r})}{\partial\mathbf{r_{3}}}\underbrace{\int_{\Omega}(\mathbf{r_{3}^{\prime}}-\mathbf{r_{3}})\nabla
Wd\mathbf{r^{\prime}}}_{C}+O(h^{2}).$
In order to ensure gradient interpolations of linear fields, it is necessary
to ensure that the discrete approximation of $A$, $B$, and $C$ are
$\displaystyle A=\begin{pmatrix}1\\\ 0\\\ 0\end{pmatrix}\quad
B=\begin{pmatrix}0\\\ 1\\\ 0\end{pmatrix}\quad C=\begin{pmatrix}0\\\ 0\\\
1\end{pmatrix}.$
By the renormalization procedure [1], we modify $\nabla W$ as follows
$\displaystyle\sum_{j}\dfrac{m_{j}}{\rho_{j}}(\mathbf{r_{j}}-\mathbf{r})\footnotesize{\bigotimes}L(\mathbf{r})\nabla
W(|\mathbf{r}-\mathbf{r}_{j}|)=\begin{pmatrix}1&0&0\\\ 0&1&0\\\
0&0&1\end{pmatrix},$
where $L(\mathbf{r})$ is a $(d,d)$ correction matrix, $d$ is the dimension of
the case. In this paper, we consider three-dimensional cases ($d=3$) and
calculate $L(\mathbf{x})$ to increase the accuracy of gradient kernel
approximation. The continuity equation is discretized by the following manner
$\dfrac{D\rho_{i}}{Dt}=\rho_{i}\sum_{j}\dfrac{m_{j}}{\rho_{j}}(\mathbf{v_{i}}-\mathbf{v_{j}})L(\mathbf{r_{i}})\nabla_{r_{i}}W(|\mathbf{r_{i}}-\mathbf{r_{j}}|,h).$
This discretized form ensures exact interpolations for both constant and
linear fields. Note that we can discretize the conservation of momentum by the
following manner
$\dfrac{D\mathbf{v_{i}}}{Dt}=-\sum_{j}m_{j}\Big{(}\dfrac{P_{i}+P_{j}}{\rho_{i}\rho_{j}}+\Pi_{ij}\Big{)}L(\mathbf{r_{i}})\nabla_{r_{i}}W(|\mathbf{r_{i}}-\mathbf{r_{j}}|,h)+\mathbf{F}.$
### 4.2 Numerical time integration with renormalization
As the other explicit hydrodynamic methods, different numerical time
integrations can be applied in SPH simulation, such as Leap-Frog, predictor-
corrector, Runge-Kutta, and Beeman schemes. The advantages of the Leap-Frog
algorithm are its low memory usage on storage and its computational
efficiency. We applied it in this paper with its improvement by using gradient
kernel renormalization. Therefore,
$\displaystyle\mathbf{r_{i}^{*}}=\mathbf{r_{i}^{n}}+\dfrac{dt}{2}\mathbf{v_{i}^{n}},$
$\displaystyle
p_{i}^{n}=\dfrac{\rho_{0}c_{0}^{2}}{\gamma}\Big{[}\Big{(}\dfrac{\rho_{i}^{n}}{\rho_{0}}\Big{)}^{\gamma}-1\Big{]},$
,
$\displaystyle\mathbf{v}_{i}^{n+1/2}=\mathbf{v}_{i}^{n-1/2}-dt\sum_{j}m_{j}\Big{(}\dfrac{p_{i}^{n}+p_{j}^{n}}{\rho_{i}^{n}\rho_{j}^{n}}+\Pi_{ij}\Big{)}L(r_{i}^{n+1/2})\nabla_{r_{i}}W(\mathbf{r}_{ij}^{n+1/2},h)+\mathbf{F}_{i}^{n+1/2},$
$\displaystyle\rho_{i}^{n+1/2}=\rho_{i}^{n-1/2}+dt\rho_{i}\sum_{j}\dfrac{m_{j}}{\rho_{j}}\big{(}\mathbf{v_{i}^{n+1/2}}-\mathbf{v_{j}^{n+1/2}}\big{)}L(r_{i}^{n+1/2})\nabla_{r_{i}}W(\mathbf{r}_{ij}^{n+1/2},h),$
$\displaystyle\mathbf{r}_{i}^{n+1}=\mathbf{r}_{i}^{n+1/2}+\dfrac{dt}{2}\mathbf{v}_{i}^{n+1/2}.$
## 5 Implementation
In the following section, the results of numerical simulations for improved
SPH method are given. This method was implemented in dam-break problem with
rigid ball structures and water waves generated by oblique piston type wave-
maker.
### 5.1 Dam-break and structure
In this implementation, we consider a rectangular tank with three-dimensional
problem, in particular on interaction between waves and structures. Here we
examine the impact of a single wave with rigid ball structures over the slope
by means of a three-dimensional SPH method. A rectangular tank contains fixed
structures and we used 10075 particles for this simulation. The geometry is
shown in Fig. 2.
Figure 2: Sketch of dam-break problem considered, side view (upper panel) and
top view (lower panel).
Fig. 3 shows the motion of a single wave which moves through rigid ball
structures in a rectangular tank. The frame at $T=0.0$ s shows the initial
configuration. In the next frame at $T=0.11175$ s, the wave generated by the
dam break and the initial layer of water on the bottom collides with the front
of the rigid ball structures and at $T=0.1425$ s the wave wraps around the
rigid ball structures. At $T=0.1725$ s, the waves collide from both sides of
the rigid ball structures then continue moving toward the right vertical wall.
The wave reflects after colliding with the opposite wall of the tank at
$T=0.3255$ s. The last movement, at $T=0.47625$ s, the reflected wave hits the
back of rigid ball structures.
Figure 3: ($T=0.0$ s) initial configuration; ($T=0.11175$ s) wave hitting the
rigid ball structures; ($T=0.1425$ s) wave wrapping around the rigid ball
structures; ($T=0.1725$ s) waves colliding after passing the rigid ball
structures; ($T=0.3255$ s) wave colliding with the opposite wall of the tank;
($T=0.47625$ s) reflected wave hitting the back of the rigid ball structures.
### 5.2 Oblique piston type wave-maker
In this implementation, we consider a rectangular tank with three-dimensional
problem, in particular the propagation of waves towards a slope. This
simulation involves a wave-maker in the form of an oscillating oblique piston
on the left-hand side and used 80682 particles. The geometry is shown in Fig.
4.
Figure 4: Sketch of oblique piston type wave-maker, side view (upper panel)
and top view (lower panel). Figure 5: ($T=0.0$ s) initial configuration;
($T=0.2865$ s) the wave is generated by the oscillating oblique piston type
wave-maker; at ($T=0.823125$ s), ($T=0.917625$ s), and ($T=0.9915$ s) the
water wave propagates onto the slope; ($T=1.18575$ s) the water wave reaches
shallow water area and hits the right wall.
Fig. 5 shows the simulation of free surface flows. The water waves generated
by oscillating piston type wave-maker were simulated. In Fig. 5 the waves are
shown propagating onto the slope. The frame at $T=0.0$ s shows the initial
configuration with water lying on the slope. In the next frame at $T=0.2865$
s, the wave is generated by the oscillating oblique piston type wave-maker. At
$T=0.823125$ s, $T=0.917625$ s, and $T=0.9915$ s the water wave propagates
onto the slope. Finally, at $T=1.18575$ s the water wave reaches shallow water
area and hits the right wall.
## 6 Summary
This paper presents the application of an improved SPH method by using
gradient kernel renormalization for simulating free surface flows. We have
implemented this method on three-dimensional cases, in particular on the
interaction between waves and structures; and the propagation of waves towards
a slope for waves generated by oblique piston type wave-maker. The three-
dimensional case of the model has been shown to produce three-dimensional
phenomenon, i.e., the collision of a single wave with rigid ball structures
and its passing around the obstacle. In summary, an improved SPH model based
on renormalization can be successfully used to simulate three-dimensional wave
problems.
## References
* [1] G. Oger, M. Doring, B. Alessandrini, P. Ferrant (2007). An improved SPH method: Towards higher order convergence. Journal of Computational Physics., 225, 1472-1492.
* [2] G. R. Liu and M. B. Liu (2003). Smoothed particle hydrodynamics: a meshfree particle method. World Scientific Publishing Co. Pte. Ltd, Singapore.
* [3] J. J. Monaghan (1988). An Introduction to SPH. Computer Physics Communications., 48, 1, 89-96.
* [4] J. J. Monaghan (1994). Simulating Free Surface Flows with SPH. Journal of Computational Physics., 110, 399 - 406.
* [5] J. J. Monaghan and J. C. Lattanzio (1985). A Refined Method for Astrophysical Problems. Astron. Astrophys., 149, 399 - 406.
* [6] J. J. Monaghan (1992). Smoothed Particle Hydrodynamics. Annu. Rev. Astron. Astrophys., 30, 543 - 574.
* [7] L. B. Lucy (1977). A numerical approach to the testing of the fission hypothesis. Astron. J., 82, 1013 - 1024.
* [8] R. A. Gingold and J. J. Monaghan (1977). Smoothed particle hydrodynamics: theory and application to non-spherical stars. Mon. Not. R. Astr. Soc., 181, 375 - 389.
|
arxiv-papers
| 2013-09-16T19:18:48 |
2024-09-04T02:49:50.972062
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Rizal Dwi Prayogo, Christian Fredy Naa",
"submitter": "Christian Naa",
"url": "https://arxiv.org/abs/1309.4074"
}
|
1309.4146
|
# Systematic Uncertainties in the Analysis of the Reactor Neutrino Anomaly
A.C. Hayes1, J.L. Friar1, G.T. Garvey1, Gerard Jungman1, G. Jonkmans2 1Los
Alamos National Laboratory, Los Alamos, NM, USA 87545
2AECL, Chalk River Laboratories, Chalk River, Ontario, Canada, K0J 1J0
###### Abstract
We examine uncertainties in the analysis of the reactor neutrino anomaly,
wherein it is suggested that only about 94% of the emitted antineutrino flux
was detected in short baseline experiments. We find that the form of the
corrections that lead to the anomaly are very uncertain for the 30% of the
flux that arises from forbidden decays. This uncertainty was estimated in four
ways, is larger than the size of the anomaly, and is unlikely to be reduced
without accurate direct measurements of the antineutrino flux. Given the
present lack of detailed knowledge of the structure of the forbidden
transitions, it is not possible to convert the measured aggregate fission beta
spectra to antineutrino spectra to the accuracy needed to infer an anomaly.
Neutrino physics conclusions based on the original anomaly need to be
revisited, as do oscillation analyses that assumed that the antineutrino flux
is known to better than approximately $4\%$.
The term “reactor neutrino anomaly” first appeared in a publication by G.
Mention et al. mention , where it generally referred to the $3\sigma$ deficit
of neutrinos detected in short-baseline reactor neutrino experiments relative
to the number predicted. The predicted number of detected neutrinos has
evolved upward over time, largely as a consequence of a predicted increase in
the energy of the neutrino flux and an increased $\bar{\nu}_{e}+p\rightarrow
n+e^{+}$ cross section associated with smaller values for the neutron
lifetime. This cross section is used to infer the neutrino flux in a
presumably well-characterized detector. The changes in the predicted neutrino
flux are mostly associated with improved knowledge of the beta decays of the
isotopes created in fission reactors. Such an anomaly would potentially be
extremely significant, if a shortfall in the detected neutrino flux could be
ascribed to $\bar{\nu}_{e}$ oscillation into a light sterile neutrino with a
mass of about 1 eV.
There is an extensive recent literature dealing with the reactor anomaly,
starting with a seminal paper by Mueller et al. mueller that reexamined the
reactor antineutrino flux. The latter publication sought to improve the
earlier flux estimates based on the ILL on-line measurements schreck ;
schreck2 ; schreck3 of the integral beta spectrum of the fission products. An
antineutrino spectrum can be inferred from a beta spectrum provided one knows
the linear combination of operators involved in the decay, the end-point
energy, and the nuclear charge. The fission beta spectra involve about 6000
beta transitions, of which about 1500 are forbidden endf-7 . Clearly some
assumptions are required in order to infer the fission antineutrino flux. The
improvements mention ; mueller on the earlier analyses of ILL integral
measurements led to an increased energy of the antineutrino flux, which was
subsequently verified in an independent analysis huber .
The present contribution examines the consequences of the forbidden
transitions known to be present (at the 30% level) in the beta decay of
fission products. We analyze the antineutrino flux, using a first-principles
derivation of the finite size (FS) and weak magnetism (WM) corrections that
were the main focus of the analyses in Refs. mueller ; huber . In addition, we
examine the shape factors needed to describe the forbidden transitions. We
find that the forbidden transitions introduce a large uncertainty in the
predicted antineutrino flux irrespective of whether the antineutrino spectra
were deduced using nuclear databases or by inverting measured aggregate
fission beta spectra. As detailed below, this finding results from the fact
that the corrections are nuclear-operator dependent and that an undetermined
combination of matrix elements contributes to non-unique forbidden
transitions.
The beta-decay spectrum $S$ for a single transition in nucleus $(Z,A)$ with
end-point energy $E_{0}=E_{e}+E_{\nu}$ is
$S(E_{e},Z,A)=S_{0}(E_{e})F(E_{e},Z,A)C(E_{e})(1+\delta(E_{e},Z,A))\,,$ (1)
where $S_{0}=G_{F}^{2}\,p_{e}E_{e}(E_{0}-E_{e})^{2}/2\pi^{3}$, $E_{e}(p_{e})$
is the electron total energy (momentum), $F(E_{e},Z,A)$ is the Fermi function
needed to account for the Coulomb interaction of the outgoing electron with
the charge of the daughter nucleus, and $C(E_{e})$ is a shape factor schopper
for forbidden transitions due to additional lepton momentum terms. For allowed
transitions $C(E)=1$. The term $\delta(E_{e},Z,A)$ represents fractional
corrections to the spectrum that were the central focus of the original
anomaly studies. The primary corrections to beta decay are radiative, finite
size, and weak magnetism, or
$\delta(E_{e},Z,A)=\delta_{\mathrm{rad}}+\delta_{\mathrm{FS}}+\delta_{\mathrm{WM}}$.
Before discussing the details of the corrections $C(E_{e})$ and
$\delta(E_{e})$, we briefly summarize the treatments used in earlier work. The
radiative corrections as derived by Sirlin sirlin were included in the
description of the beta spectra (though not in the antineutrino spectra) in
the original analyses of Schreckenbach _et al._ schreck ; schreck2 ; schreck3
. In the later ILL work schreck2 ; schreck3 an approximation for the FS and
WM corrections was included by first deducing the antineutrino spectrum from
the measured beta spectra without these corrections, and then applying a
linear correction to the deduced antineutrino spectrum of the form,
$\delta_{\mathrm{FS}}+\delta_{\mathrm{WM}}=0.0065(E_{\nu}-{\rm 4\;MeV})$. In
that work no corrections were made for the shape factors $C(E_{e})$. In the
analyses of Refs. mueller ; huber ; fallot an approximation (derived by Vogel
vogel-2 ) for the FS and WM corrections was applied on a transition-by-
transition basis. In Refs. mueller ; fallot the shape factor appropriate for
unique forbidden transitions was used for all forbidden transitions. In Ref.
huber it was argued that these shape factors only play a small role in
inferring antineutrino spectra from measured beta spectra because for
vanishing electron mass, $m_{e}$, they are symmetric under
$E_{e}\leftrightarrow E_{\nu}$. In the present work, we derived ab initio
analytic expressions for the FS and WM corrections for allowed GT transitions,
as well as WM and shape factors for first-forbidden GT operators. We used the
radiative corrections derived by Sirlin sirlin .
We now turn to the form of the corrections. The attractive Coulomb interaction
increases the electron density near the nuclear surface and increases the
beta-decay rate, while the finite nuclear size decreases the electron density
and decreases the rate (relative to the point-nucleus Fermi function). Using
first-order perturbation theory in $Z\alpha$, we find that the finite-size
correction to the Fermi function, $\delta_{\mathrm{FS}}$, for allowed GT
transitions is
$\delta_{\mathrm{FS}}=-\frac{3}{2}\frac{Z\alpha}{\hbar
c}\left<r\right>_{(2)}\left(E_{e}-\frac{E_{\nu}}{27}+\frac{m_{e}^{2}c^{4}}{3E_{e}}\right)\\\
\;\;.$ (2)
The quantity $\langle r\rangle_{(2)}=\int d^{3}r\,\rho_{W}(r)\int
d^{3}s\,\rho_{\rm ch}(s)\mid\vec{r}-\vec{s}\mid$ is the first moment of the
convoluted nuclear weak and charge densities (called a Zemach moment zemach ).
We assume uniform distributions of radius $R$ for the weak and charge
densities, for which $\left<r\right>_{(2)}=\frac{36}{35}R$ friar1 . The FS
corrections do depend on the beta-decay operator, but in this work we always
use Eq. (2) and $R=1.2\;A^{1/3}$ fm.
The WM correction arises from the interference of the magnetic moment
distribution of the vector current, $\vec{J}_{V}=\vec{\nabla}\times\vec{\mu}$,
with the spin distribution $\vec{\Sigma}$ of the axial current. We derived the
WM corrections for allowed and first-forbidden operators. There are four
possible operators in the case of first-forbidden GT transitions, and all have
well-defined WM corrections, as listed in Table 1.
Table 1: The shape factors and leading-order weak magnetism corrections to allowed and first-forbidden Gamow-Teller beta decays are shown in the top panel. The shape factors for allowed and first-forbidden Fermi beta decays are shown in the bottom panel. All agree with Ref. milliner for $Z=0$. The entries for $\vec{J}_{V}$ and $\rho_{A}$ are discussed in siegert . The weak magnetism correction for $\vec{J}_{V}$ involves the unknown overlap of very different $1^{-}$ matrix elements and is therefore not listed. The nucleon isovector magnetic moment is $\mu_{v}=4.7$, $M_{N}$ is the nucleon mass, $g_{A}$ is the axial vector coupling constant, and $\beta=p_{e}/E_{e}$. No meson currents were used in the magnetic moment operator, and a truncated orbital current led to the factor of “1/2” in $\delta_{\mathrm{WM}}$. Classification | $\;\Delta J^{\pi}\;$ | Operator | Shape Factor $C(E_{e})$ | Fractional Weak Magnetism Correction $\delta_{\mathrm{WM}}(E_{e})$
---|---|---|---|---
Allowed GT | $1^{+}$ | $\Sigma\equiv\sigma\tau$ | 1 | $\frac{2}{3}\left[\frac{\mu_{v}-1/2}{M_{N}g_{A}}\right](E_{e}\beta^{2}-E_{\nu})$
Non-unique 1st Forbidden GT | $0^{-}$ | $\left[\Sigma,r\right]^{0-}$ | $p_{e}^{2}+E_{\nu}^{2}+2\beta^{2}E_{\nu}E_{e}$ | 0
Non-unique 1st Forbidden $\rho_{A}$ | $0^{-}$ | $\left[\Sigma,r\right]^{0-}$ | $\lambda\,E_{0}^{2}$ | 0
Non-unique 1st Forbidden GT | $1^{-}$ | $\left[\Sigma,r\right]^{1-}$ | $p_{e}^{2}+E_{\nu}^{2}-\frac{4}{3}\beta^{2}E_{\nu}E_{e}$ | $\;\;\;\;\left[\frac{\mu_{v}-1/2}{M_{N}g_{A}}\right]\left[\frac{(p_{e}^{2}+E_{\nu}^{2})(\beta^{2}E_{e}-E\nu)+2\beta^{2}E_{e}E_{\nu}(E_{\nu}-E_{e})/3}{(p_{e}^{2}+E_{\nu}^{2}-4\beta^{2}E_{\nu}E_{e}/3)}\right]$
Unique 1st Forbidden GT | $2^{-}$ | $\left[\Sigma,r\right]^{2-}$ | $p_{e}^{2}+E_{\nu}^{2}$ | $\frac{3}{5}\left[\frac{\mu_{v}-1/2}{M_{N}g_{A}}\right]\left[\frac{(p_{e}^{2}+E_{\nu}^{2})(\beta^{2}E_{e}-E\nu)+2\beta^{2}E_{e}E_{\nu}(E_{\nu}-E_{e})/3}{(p_{e}^{2}+E_{\nu}^{2})}\right]$
Allowed F | $0^{+}$ | $\tau$ | 1 | 0
Non-unique 1st Forbidden F | $1^{-}$ | $r\tau$ | $p_{e}^{2}+E_{\nu}^{2}+\frac{2}{3}\beta^{2}E_{\nu}E_{e}$ | 0
Non-unique 1st Forbidden $\vec{J}_{V}$ | $1^{-}$ | $r\tau$ | $E_{0}^{2}$ | -
Our FS and WM corrections for allowed GT transitions are identical to those
derived by Holstein holstein , and differ from the forms used in mueller ;
huber ; fallot ; vogel-2 . The first-forbidden shape factors, $C(E_{e})$,
which depend on the operator in question, were derived and are displayed in
Table 1. To determine the implications of the corrections we take advantage of
the fact that the ENDF/B-VII.1 decay library endf-7 contains branching ratios
and end-point energies for over 90% of the transitions making up the full
spectra. Thus, the library allows a representative estimate of the effect of
the corrections, provided we know what assumptions to make about the forbidden
transitions.
We first fit the Schreckenbachschreck2 electron spectrum for 235U before
examining the effect of the corrections in an analysis based on the summation
of the known beta branches of the fission products. We include the corrections
both with and without a treatment of the forbidden transitions. We fit the
beta spectrum with 40 fictitious equally spaced end-point energies using non-
negative least-squares fitting nnls . The charge Z associated with each end-
point is assigned according to the method given in schreck2 . There is no
unique physical prescription for beta-decay operator assignments to the
fictitious end-points. For this reason we examine four prescriptions: (1) all
transitions are assumed to be allowed; (2) all end-point energies can be
associated with either an allowed or forbidden transition; (3) 30% of the
branches are selected to be forbidden at equal energy intervals; (4) 30% of
the branches are selected to be forbidden with a bias towards higher energies.
In addition, we examine fits in which the operator determining the forbidden
decays was taken to be $[\Sigma,r]^{0-},[\Sigma,r]^{1-}$, $[\Sigma,r]^{2-}$ or
a combination of these. As a natural consequence of the non-negative least-
squares procedure, the fit results in a significant fraction of the end-points
having zero amplitude. The subset of end-points with non-zero amplitudes
varies depending on the operator assignment.
We find excellent fits to the electron spectrum in all cases. However,
different treatments of the forbidden transitions can lead to antineutrino
spectra that differ both in shape and magnitude at about the 4% level. Two
examples are shown in Fig.1, where we present the fits obtained when the WM
and FS corrections are included. In one case all transitions are assumed to be
allowed, while in the second case the best fit results from about 25%
forbidden decays. For the assumption of all allowed transitions, we see a
systematic increase of about 2.2% in the number of antineutrinos relative to
Schreckenbach, while including forbidden transitions leads to no increase
relative to Schreckenbach. Other prescriptions for the forbidden transitions
lead to changes relative to the Schreckenbach antineutrino spectrum ranging
from 0-4%, including changes in the shape of the antineutrino spectrum. These
examples help to clarify the difficulty in inferring with high certainty the
antineutrino spectrum from a measured electron spectrum when information on
the forbidden transitions is not available.
Figure 1: The fit to the electron spectrum for 235U (left) for two different
assumptions on how to treat forbidden transitions, and the ratio of the
corresponding antineutrino spectra to that of Schreckenbach (right). The
electron data are those of schreck2 . The electron spectra are fit assuming
(a) all allowed GT branches, or (b) up to 30% forbidden GT transitions. In
both cases the WM and FS corrections are included. When folded over the
neutrino detection cross section, the case for all allowed (25% forbidden)
transitions results in a 2.2% (0.06%) increase in the number of detectable
antineutrinos.
In calculating the aggregate fission spectra from the database library an
analogous uncertainty arises because detailed structure information does not
exist for the majority of the roughly 1500 forbidden transitions. In addition,
several of the transitions are forbidden at second or higher order, for which
analytic corrections are not available. Thus, there is no clear prescription
for applying the corrections to this component of the spectra, and it is more
beneficial to examine the effect of different approximations in order to
estimate the uncertainty involved. In all approximations we treat unique
forbidden transitions as unique first-forbidden GT transitions, and treat non-
unique forbidden transitions in one of the following ways: (1) as allowed GT;
(2) as unique first-forbidden GT with the operator
$\left[\Sigma,r\right]^{2-}$; (3) with the operator
$\left[\Sigma,r\right]^{0-}$; (4) with the operator
$\left[\Sigma,r\right]^{1-}$. None of the these treatments is correct, but
they provide estimates for changes in the spectra induced by forbidden
transitions.
The aggregate fission beta spectrum under equilibrium reactor burning
conditions for a given actinide is determined by the beta spectra
$S(E_{e},Z_{i},A_{i})$ of the individual unstable fission fragments weighted
by their cumulative fission yields, $Y_{F_{i}}$ tal :
$N_{\beta}(E_{e})=\sum_{F_{i}}Y_{F_{i}}S(E_{e},Z_{i},A_{i}).$ (3)
The beta spectrum $S$ for each fragment $(Z_{i},A_{i})$ summed over all decay
branches must be normalized to unity: $\int S(E,Z,A)\,dE=1$. Thus, Eq. (3) is
a statement that under equilibrium burning conditions the beta-decay rates are
determined by the fission rate nonequilibrium . If the antineutrino spectrum
is inferred from a measured aggregate beta spectrum, Eq. (3) must be replaced
by a sum over a set of end-point energies $\\{E_{0_{i}}\\}$, weighted by a
fitted set of coefficients $\\{a_{i}\\}$:
$N_{\beta}=\sum_{i}a_{i}S(E_{e},E_{0_{i}})$.
There is no unique method for determining the uncertainty in the antineutrino
spectrum introduced by the forbidden transitions. Another possibility (in
addition to the fitting exercise above) is to consider changes in the bi-
variant function $k(E_{e},E_{\nu})$, where
$k(E_{e},E_{\nu})=N_{\nu}(E_{\nu})/N_{\beta}(E_{e})$. If $k(E_{e},E_{\nu})$
only changes by some small percentage for some path in the $(E_{e},E_{\nu})$
plane as we change our treatment of forbidden transitions, then there exists a
prescription for inferring the antineutrino spectrum to that accuracy. We
calculated the function $k(E_{e},E_{\nu})$ for each of our four assumptions
(above). We found no path in the $(E_{e},E_{\nu})$ plane that left
$k(E_{e},E_{\nu})$ unchanged by as little as 5% as our assumptions for the
forbidden transitions changed. Fig. (2) depicts the result for the path
$E_{\nu}=K_{\beta}\equiv E_{e}-m_{e}c^{2}$ borovoi . Similar or larger
differences were found for all other paths. The very non-smooth and non-linear
shape of $k$ arises from the shape factors $C(E_{e})$.
Figure 2: The ratio of the function $k(E_{e},E_{\nu})$ for 235U using
ENDFB/VII.1 relative to using $E_{\nu}=K_{\beta}$ borovoi .
A third prescription for estimating the uncertainties is to examine the rate
of change in the antineutrino spectrum relative to the rate of change in the
beta spectrum, using the fact that the beta spectrum is fit to amplitudes
${a_{i}}$ on a fixed grid of end-point energies ${E_{0_{i}}}$. We calculated
$T=\sum_{i}[\partial N_{\nu}(E_{\nu})/\partial a_{i}]/[\partial
N_{\beta}(E_{e})/\partial a_{i}]$, and examined the changes in $T$ as the
assumptions for the forbidden transitions were varied, and again found no path
in $(E_{e},E_{\nu})$ over which the changes were $<$ 5%.
Figure 3: Different treatments of the forbidden GT transitions contributing to
the antineutrino spectrum summed over all actinides in the fission burn in
mid-cycle kopeikin of a typical reactor. The left panel shows the ratio of
these antineutrino spectra relative to that using the assumptions of Ref.
schreck2 . The right panel shows the spectra weighted by the detection cross
section, where the additional curve in black uses the assumptions of Ref.
schreck2 . The spectra are strongly distorted by the forbidden operators,
being lower below the peak and in some cases more than 20% larger above the
peak than Ref. schreck2 . The corresponding change in the number of detectable
antineutrinos relative to schreck2 is -0.75%, 5.8% and 1.85% for the
$0^{-},1^{-}$, and $2^{-}$ forbidden operators, respectively.
Our final method considers the ratio of the actual antineutrino spectra
themselves in Fig. (3), although this method does not take into account the
corresponding changes in the beta spectra. We carried out identical analyses
of the role of the corrections and the associated uncertainties for the other
actinides 239,241Pu and 238U, and found very similar results. Fig. (3) shows
the change in the total antineutrino spectrum at a representative time
kopeikin in mid-cycle in the reactor burn history. The antineutrino spectra
differ significantly, depending on our treatment of the forbidden transitions.
The cross-section-weighted spectra are quite distorted, being lower than the
Schreckenbach-Vogel schreck2 ; vogel-1 ; explain spectra up to the peak, and
higher or lower above the peak depending on the operator. The actual spectrum
is unlikely to be as distorted as in Fig. (3) because no single operator
dominates the forbidden transitions.
The original ILL analysis schreck2 assumed that the $Z$ of the daughter
fragments satisfies $Z_{i}=49.5-0.7E_{0_{i}}-0.09E_{0_{i}}^{2}$ for $Z_{i}\geq
34$. We find that taking the $Z$ of the fission fragments directly from
ENDF/B-VII.1 instead would increase the antineutrino spectrum by less than 1%
for $E_{\nu}\leq$ 7 MeV, and less than 1.5% for $E_{\nu}>$ 7 MeV.
In summary we find that the component of the aggregate fission spectra
containing approximately 30% forbidden transitions introduces a large
uncertainty (about 4%) in the predicted shape of the antineutrino flux emitted
from reactors. We have examined the uncertainties in four different ways. If
all forbidden transitions are treated as allowed GT transitions, the
antineutrino spectra are systematically increased above 2 MeV, as was the
conclusion in the earlier papers on the anomaly. However, when the forbidden
transitions are treated in various approximations, the shape and magnitude of
the spectra are changed significantly. Earlier analyses only looked at one
prescription for these transitions. The uncertainty introduced by our lack of
knowledge on how to treat these transitions is as large as the size of the
anomaly. It should also be noted that there are additional uncertainties due
to our oversimplified treatment of the allowed weak magnetism operator (viz.,
no meson currents and a truncated orbital current), as well as a simplified
finite-size treatment for forbidden transitions. These considerations also
apply to fission antineutrino fluxes in medium and long baseline reactor
experiments, implying a 4% uncertainty on the fission antineutrino flux in
those experiments. Reducing the uncertainty within a purely theoretical
framework would be difficult. An improvement will require either direct
measurements of the antineutrino flux or a substantial improvement in our
knowledge of the dominant forbidden beta transitions.
## References
* (1) G. Mention, et al., Phys. Rev. D 83 073006 (2011).
* (2) Th. A. Mueller, et al., Phys. Rev. C 83, 054615 (2011).
* (3) K. Schreckenbach, H. R. Faust, F. von Feilitzsch, A. A. Hahn, K. Hawerkamp, and J. L. Vuilleumier, Phys. Lett. 99B, 251 (1981); F. von Feilitzsch, A. A. Hahn, and K. Schreckenbach, Phys. Lett. 118B, 162 (1982).
* (4) K. Schreckenbach, G. Colvin, W. Gelletly, and F. von Feilitzsch, Phys. Lett. 160B 325 (1985). The data from Table 1 and the antineutrino predictions of this reference will be denoted above simply as “Schreckenbach.”
* (5) A. A. Hahn, K. Schreckenbach, W. Gelletly, F. von Feilitzsch, G. Colvin, and B. Krusche, Phys. Lett. B218, 365 (1989).
* (6) M. B. Chadwick, Nucl. Data Sheets 112, 2887 (2011).
* (7) P. Huber, Phys. Rev. C 84, 024617 (2011).
* (8) H. F. Schopper, Weak Interactions and Nuclear Beta Decay, North-Holland, Amsterdam, 1966.
* (9) A. Sirlin, Phys. Rev. 164, 1767 (1967); A. Sirlin, Phys. Rev. D 84, 014021 (2011).
* (10) M. Fallot, S. Cormon, M. Estienne, A. Algora, V. M. Bui, et al., Phys. Rev. Lett. 109, 202504 (2012).
* (11) P. Vogel, Phys. Rev. D 29, 1918 (1984) and private communication.
* (12) C. Zemach, Phys. Rev. 104, 1771 (1956).
* (13) J. L. Friar, Ann. Phys. (N.Y.) 122, 151 (1979).
* (14) B. R. Holstein, Phys. Rev. C 9, 1742 (1974); J. L. Friar and I. Sick, Phys. Lett. B579, 285 (2004). Equation (2) of the latter allows conversion between the two results.
* (15) C. Lawson and R. Hanson, _Solving Least Squares Problems_ , Prentice-Hall, Englewood Cliffs, New Jersey, 1974.
* (16) D. J. Millener, D. E. Alburger, E. K. Warburton, and D. H. Wilkinson, Phys. Rev C 26, 1167 (1982).
* (17) Nuclear operators $\vec{J}_{V}$ and $\rho_{A}$ are proportional to a nucleon velocity ($p/M_{N}$), have negative parity, and are therefore first-forbidden. Replacing them by the analogous operators proportional to $E_{0}\,r$ requires only CVC for $\vec{J}_{V}$, but is an approximation for $\rho_{A}$ made exact by introducing an unknown dimensionless nuclear-structure parameter $\lambda$ that differs from case to case.
* (18) T. R. England and B. F. Rider: “Evaluation and Compilation of Fission Product Yields,” LA-UR-94-3106 (1994).
* (19) Off-equilibrium corrections are needed for fission fragments that decay with half-lives long compared to the timescales involved in the experiments in schreck ; schreck2 ; schreck3 .
* (20) A. A. Borovoi, V. I. Kopeikin, L. A. Mikaelyan, and S. V. Tolokonnikov, Sov. Jour. Nucl. Phys. 36, 232 (1982).
* (21) V. I. Kopeikin, Phys. Atomic Nuclei 75 143 2012. The mid-cycle combination of burning actinides is 56% 235U, 31% 239Pu, 7% 238U and 6% 214Pu.
* (22) P. Vogel, G. K. Schenter, F. M. Mann, and R. E. Schenter, Phys. Rev. C 24, 1543 (1981).
* (23) “Schreckenbach-Vogel” means that we used the spectra for 235U, 239Pu, and 241Pu from schreck2 and 238U from vogel-1 .
|
arxiv-papers
| 2013-09-17T01:07:29 |
2024-09-04T02:49:50.982733
|
{
"license": "Public Domain",
"authors": "A.C. Hayes, J.L. Friar, G.T. Garvey, G. Jungman, Guy Jonkmans",
"submitter": "Jim Friar",
"url": "https://arxiv.org/abs/1309.4146"
}
|
1309.4234
|
Three-Dimensional Smoothed Particle Hydrodynamics Simulation
for Liquid Metal Solidification Process
Raden Ahnaf Faqih S.a,b, Christian Fredy Naab
aGraduate School of Natural Science and Technology, Kanazawa University,
Kakuma, Kanazawa 920-1192 Japan,
bFaculty of Mathematics and Natural Sciences, Institut Teknologi Bandung, Jl.
Ganesha 10, Bandung 40132 Indonesia,
E-mail: [email protected], [email protected]
Abstract.The solidification behavior of liquid metal in a container under
rapid cooling process is one of the major concerns to be analyzed. In order to
analyze its fundamental behavior, a three-dimensional (3D) fluid dynamics
simulation was developed using a particle-based method, known as the smoothed
particle hydrodynamics (SPH) method. Governing equations that determine the
fluid motion and heat transfer involving phase change process are solved by
discretizing their gradient and Laplacian term with the moving particles and
calculating the interaction with its neighboring particles. The results
demonstrate that the SPH mehod can successfully reproduce the behavior and
defect prediction of liquid metal solidification process.
Keywords: liquid metal, particle-based method, rapid cooling, solidification,
SPH
## 1 Introduction
Understanding the solidification behavior of liquid metal is one of the
important things for industrial manufacturing field. Sometimes several defects
occur if the metal processing is based on the solidification process such as
disperse shrinkage, shrinkage cavity, blind cavity, linear cracks, gas
porosity, and so on. Various studies are performed for understanding these
behaviors. Nowadays, many researchers try to understand it by using numerical
study. They make a mathematical model for the multiphase flow, heat transfer,
phase transition, and then demonstrate it through a numerical simulation. A
review of mathematical modelling of solidification process is discussed in
[1]. As to the numerical simulation, it can be divided into two methods, the
mesh-based method and particle-based method. Conventional mesh-based methods,
in general, encounter difficulties in representing multiphase flow. To
overcome this problem, the particle-based method which does not need to
generate computational mesh can be used. Several particle-based methods have
been developed in recent years. One of the particle-based methods which can be
applied to simulate incompressible multiphase flows is smoothed particle
hydrodynamics. SPH was originally introduced by Lucy [2] and Gingold &
Monaghan [3] to simulate three-dimensional astrophysics problems. However,
this method was eventually developed for simulating fluid dynamics and heat
transfer problems. SPH method applied to heat transfer for solidification can
be seen in [4]. They performed the solidification process for water. In other
case, Cleary et al [5] have simulated liquid metal solidification under high
pressure die casting. However, they did not explain in detail about the
interaction between liquid and solidified metal.
The objectives of this paper are to understand and analyze the application of
SPH method for liquid metal solidification. We also want to analyze the
behavior and defect prediction of liquid metal especially for the case of
rapid solidification process. We first apply the SPH formulation for the fluid
flow by attaching a temperature dependence of viscosity. We then apply the SPH
method for heat transfer which is based on the enthalpy formulation. To verify
that the heat transfer works well, we test the heat transfer model by using 1D
static particle. For the phase transition model, we consider the nonisothermal
phase change problem. We also consider the interaction between different
phases and examine it for 3D simulation case.
## 2 Numerical model
### 2.1 The SPH method
In SPH, the fluid is represented by a set of particles that can move freely.
Each particle carries some fundamental physical properties, such as mass,
position, velocity, density, and any other related properties. The value of
any function $f$ at a particle can be approximated by summing over the
properties of its neighboring particles. The SPH interpolation of the function
$f$ of particle $i$ at position $\textbf{r}_{i}$ is approximated by:
$f(\textbf{r}_{i})\approx\sum_{j}\frac{m_{j}}{\rho_{j}}f_{j}W(\textbf{r}_{i}-\textbf{r}_{j},h),$
(1)
where index $j$ corresponds to any neighboring particle of particle $i$,
$m_{j}$ and $\rho_{j}$ are the mass and the density for particle $j$, $f_{j}$
is the value of $f$ for particle $j$, the function $W$ is an interpolation
smoothing kernel, and $h$ is a smoothing length that defines the radius of
influence around the particle $i$.
There are many kinds of smoothing kernel functions in the literature and here
we choose cubic spline kernel function as below [6]:
$W_{ij}=W(\textbf{r}_{ij},h)=C\begin{cases}\frac{2}{3}-\left(\frac{|\textbf{r}_{ij}|}{h}\right)^{2}+\frac{1}{2}\left(\frac{|\textbf{r}_{ij}|}{h}\right)^{3}&\text{,
}0\leq\frac{|\textbf{r}_{ij}|}{h}<1\\\
\frac{1}{6}\left(2-\left(\frac{|\textbf{r}_{ij}|}{h}\right)\right)^{3}&\text{,
}1\leq\frac{|\textbf{r}_{ij}|}{h}<2\\\ 0&\text{,
}2\geq\frac{|\textbf{r}_{ij}|}{h}\end{cases}$ (2)
where $\textbf{r}_{ij}=\textbf{r}_{i}-\textbf{r}_{j}$ and $C=3/(2\pi h^{2})$
for 3D simulation. The radius of influence of this function is $2h$. It means
that the particle $i$ will not be affected by any neighboring particles
farther than $2h$. However, the closer distance of the neighboring particles
to the particle $i$, the greater will be the influence accepted by the
particle $i$.
The derivative of the function $f$ is obtained by differentiating the
interpolation Eq.(1) which is given by:
$\nabla
f(\textbf{r}_{i})\approx\sum_{j}\frac{m_{j}}{\rho_{j}}f_{j}\nabla_{i}W_{ij}$
(3)
where
$\nabla_{i}W_{ij}=\nabla
W(\textbf{r}_{ij},h)=\frac{1}{h}\frac{\textbf{r}_{ij}}{|\textbf{r}_{ij}|}F(|\textbf{r}_{ij}|,h)$
and $F$ is the derivative of the kernel function Eq.(2).
### 2.2 Fluid model
The governing equations of fluids in SPH method are based on the Navier-Stokes
equations in the Lagrangian form. The basic equations are given by [7]:
$\displaystyle\frac{d\rho}{dt}$ $\displaystyle=$
$\displaystyle-\rho\nabla\cdot\textbf{v}$ (4)
$\displaystyle\rho\frac{d\textbf{v}}{dt}$ $\displaystyle=$
$\displaystyle-\nabla p+\mu\nabla^{2}\textbf{v}+\textbf{F}$ (5)
where $t$ is time, v is velocity, $p$ is pressure, $\mu$ is viscosity, and F
is external force. Eq.(4) is known as the continuity equation which describes
the evolution of the fluid density over time. Eq.(5) is known as the momentum
equation which describes the acceleration of the fluid. The first term on the
right hand side involving the pressure gradient corresponds to the pressure
force and the second term involving viscosities corresponds to the viscous
force. The external forces used here are gravitational force and repulsive
force from the boundaries.
By employing the SPH interpolation of Eq.(3) to the Eq.(4) and Eq.(5) then the
SPH representation of the continuity and the momentum equation can be written
as [6]:
$\displaystyle\frac{d\rho_{i}}{dt}$ $\displaystyle=$
$\displaystyle\sum_{j}m_{j}(\textbf{v}_{i}-\textbf{v}_{j})\cdot\nabla W_{ij},$
(6) $\displaystyle\frac{d\textbf{v}_{i}}{dt}$ $\displaystyle=$
$\displaystyle-\sum_{j}m_{j}\left[\left(\frac{P_{i}}{\rho_{i}^{2}}+\frac{P_{j}}{\rho_{j}^{2}}\right)-\frac{\xi}{\rho_{i}\rho_{j}}\frac{4\mu_{i}\mu_{j}}{\left(\mu_{i}+\mu_{j}\right)}\frac{\textbf{v}_{ij}\textbf{r}_{ij}}{\textbf{r}_{ij}^{2}+\eta^{2}}\right]\nabla_{i}W_{ij}+\textbf{F}.$
(7)
The viscous force used here is the laminar viscosity term which was introduced
by Morris and Monaghan [8] with $\mu_{i}$ and $\mu_{j}$ the dynamic viscosity
of the fluid corresponding to particles $i$ and $j$. Further, $\xi=4.963$ is a
viscous scaling factor [9] and $\eta$ is a small parameter to prevent
singularity when $\textbf{r}_{ij}$ goes to zero (usually $\eta\sim 0.1h$).
An equation of state is required to calculate the pressure in Eq.(7). The
equation of state used here is quasi-compressible form which is calculated by
using the density information from Eq.(6). The equation of state is given by
[9]:
$P=\beta\left[\left(\frac{\rho}{\rho_{0}}\right)^{\gamma}-1\right]$ (8)
where $\beta=c^{2}\rho_{0}/\gamma$ is the magnitude of the pressure,
$\rho_{0}$ is the reference density, $\gamma=7$ for liquid metal, and $c$ is
the speed of sound. The value of speed of sound should be large enough to
ensure that the density fluctuation is less than 1% or close to the
incompressible flow.
### 2.3 Heat transfer and phase transition model
The heat transfer process occurs by cooling down the system through the
boundaries. Here, the boundaries are modeled by a set of SPH particles which
have a certain temperature and thermal conductivity. The temperature and
thermal conductivity of the boundary particles are assumed to be constant
during the solidification process. We also assume the boundary as adiabatic
boundary where there is no flux to the ambient environment. The heat exchange
between the particles (liquid, solid, and boundary) occurs by considering the
heat conduction. The model of heat conduction is based on the enthalpy method
which is given by:
$\frac{dH}{dt}=\frac{1}{\rho}\nabla(k\nabla T)$
where $H$ is enthalpy, $k$ is thermal conductivity, and $T$ is temperature.
The SPH formulation of this equation is approximated by using the modified SPH
approximation for second derivative which is developed by Cleary and Monaghan
[10]:
$\frac{dH_{i}}{dt}=\sum_{j}\frac{m_{j}}{\rho_{i}\rho_{j}}\frac{4k_{i}k_{j}}{(k_{i}+k_{j})}(T_{i}-T_{j})\frac{\textbf{r}_{ij}\cdot\nabla_{i}W_{ij}}{\textbf{r}_{ij}^{2}+\eta^{2}}.$
(9)
This equation ensures that the heat flux is automatically continuous across
the different material interfaces, such as between the mold and liquid metal.
This also allows multiple phases with different conductivities to be
accurately simulated [10].
The temperature in Eq.(9) is calculated by considering the relationship
between enthalpy and temperature. In the liquid metal solidification case, the
phase change occurs over a range of temperature. The phase transition goes
through an intermediate phase which is known as mushy phase. So, the
nonisothermal phase change is involved to correlate the enthalpy and
temperature. The illustration of the nonisothermal phase change is shown in
Figure 1. This relation ensures that the temperature is continuous across the
different phases.
Figure 1: Relationship between enthalpy and temperature for nonisothermal
phase change
The relationship between enthalpy and temperature can be written as [11]:
$H=\begin{cases}\int_{T_{ref}}^{T}c_{s}(T)dT&\text{, }T\leq T_{s}\text{
$~{}~{}~{}~{}~{}~{}~{}$ (Solid)},\\\
\int_{T_{ref}}^{T_{s}}c_{s}(T)dT+\int_{T_{s}}^{T}\frac{\partial L}{\partial
T}dT&\text{, }T_{s}<T\leq T_{l}\text{$~{}~{}$(Mushy)},\\\
\int_{T_{ref}}^{T_{s}}c_{s}(T)dT+L+\int_{T_{l}}^{T}c_{l}(T)dT&\text{,
}T>T_{l}\text{$~{}~{}~{}~{}~{}~{}~{}~{}$ (Liquid)},\end{cases}$ (10)
where $T_{ref}$ is the reference temperature, $T_{l}$ and $T_{s}$ are
respectively the solidus and liquidus temperatures, and $c_{l}$ and $c_{s}$
are the specific heats in the solid and liquid phases, respectively. If we
consider the specific heats to be constant then Eq.(10) becomes:
$H=\begin{cases}c_{s}T&\text{, }T\leq T_{s},\\\
c_{s}T_{s}+c_{sl}(T-T_{s})&\text{, }T_{s}<T\leq T_{l},\\\
c_{s}T_{s}+c_{sl}(T_{l}-T_{s})+c_{l}(T-T_{l})&\text{, }T>T_{l},\end{cases}$
The temperature of each particle is then calculated by using this relation in
the temperature term:
$T_{i}=\begin{cases}\frac{H_{i}}{c_{s}}&\text{, }H_{i}\leq
H_{s}=c_{s}T_{s}\text{
$~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}$
(Solid)},\\\ T_{s}+\frac{H_{i}-H_{s}}{c_{sl}}&\text{, }H_{s}<H_{i}\leq
H_{l}=H_{s}+c_{sl}(T_{l}-T_{s})\text{$~{}$(Mushy)},\\\
T_{s}+\frac{H_{i}-H_{l}}{c_{l}}&\text{,
}H_{i}>H_{l}\text{$~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}$
(Liquid)},\end{cases}$
where $c_{sl}=L/(T_{l}-T_{s})$. The latent heat $L$ describes the energy
released by a particle to change the phase from liquid to solid.
### 2.4 Liquid-solid interaction model
The phase transition of a particle from liquid to solid occurs when the
particle temperature is below the solidus temperature. Here, we modelled the
solid particles as a viscous pseudo-fluid. It means that the behavior of solid
particles is like a fluid (liquid or mushy) particles but they move under high
viscosity. This approach is used to keep the forces acting on the solid
particles and fluid particles maintained. If the solid particle has more than
two of its neighbors being solid phase then they act and move together as a
solid group (rigid body). The illustration how to create the solid group is
shown in Figure 2. First, each solid particle is assigned an initial index, as
shown in Figure 2. After that, each particle index is updated by the maximum
index of its neighbor particle (Figure 2). This process is repeated until each
solid particle index converges to the maximum index in each group (Figure 2).
After some steps, we will only have the maximum index on each group and it
means the number of the solid group is determined by how many remaining
indices we have. For example in the Figure 2, we have two separate solid
groups because finally we have two indices that are 13 and 20.
Figure 2: Illustration of calculating the number of solid groups
Here, we consider the solid group as an undeformable object. In order to
prevent the deformation, the solid particle’s movement should be converted to
the translational and rotational motion related to the center of mass of its
solid group. The center of mass is calculated by using the following equation:
$\textbf{r}_{c}=\frac{1}{N}\sum_{j}\textbf{r}_{j}$
where $N$ is the number of solid particle in the solid group.
The translational and rotational motion is considered by calculating the
moment of inertia and angular velocity of the solid group. The moment of
inertia, $I$, and angular velocity, $\omega$, is calculated by the following
equations:
$I=\sum_{i}|\textbf{q}_{i}|^{2}\;\text{and
}\omega=\frac{1}{I}\sum_{i}\textbf{q}_{i}\times\textbf{v}_{i}.$
Here, $\textbf{q}_{i}=\textbf{r}_{i}-\textbf{r}_{c}$ is the relative position
of particle $i$ to the center of mass.
During the simulation, the center of mass and the rotation matrix are updated
by using the following equations:
$\textbf{r}_{c}(t+\Delta t)=\textbf{r}_{c}(t)+\Delta
t\textbf{v}_{c}\left(t+0.5\Delta t\right)\;\text{and }\textbf{R}(t+\Delta
t)=\textbf{R}(t)+\Delta t\textbf{Z}\left(t+0.5\Delta t\right)$
where $\textbf{v}_{p}$ is the center of mass velocity and the angular velocity
tensor $\textbf{Z}(t)$ in is given by:
$\textbf{Z}(t)=\left[\begin{array}[]{ccc}0&-\omega_{z}&\omega_{y}\\\
\omega_{z}&0&-\omega_{x}\\\ -\omega_{y}&\omega_{x}&0\\\
\end{array}\right]\textbf{R}(t).$
The rotation matrix in 3D coordinates system is given by:
$\textbf{R}_{x}(\theta)=\left[\begin{array}[]{ccc}1&0&0\\\
0&\cos\theta&\sin\theta\\\ 0&-\sin\theta&\cos\theta\\\
\end{array}\right],\;\textbf{R}_{y}(\theta)=\left[\begin{array}[]{ccc}\cos\theta&0&-\sin\theta\\\
0&1&0\\\ \sin\theta&\omega_{x}&\cos\theta\\\
\end{array}\right],\;\textbf{R}_{z}(\theta)=\left[\begin{array}[]{ccc}\cos\theta&\sin\theta&0\\\
-\sin\theta&\cos\theta&0\\\ 0&0&1\\\ \end{array}\right].$
After we get the angular velocity and the center of mass, the relative
velocity and position of each solid particle to the center of mass can be
written as below:
$\displaystyle\textbf{v}_{i}(t+\Delta t)$ $\displaystyle=$
$\displaystyle\textbf{v}_{c}(t)+\omega\times\textbf{q}_{i}(t),$
$\displaystyle\textbf{r}_{i}(t+\Delta t)$ $\displaystyle=$
$\displaystyle\textbf{R}(t+\Delta
t)\left[\textbf{r}_{i}(t)-\textbf{r}_{c}(t)\right]+\textbf{r}_{c}(t+\Delta
t).$
## 3 Numerical simulations
### 3.1 One-dimensional sinusoidal temperature
We simulate the cooling down process in 1D case to verify that the SPH method
for heat conduction works well. We use 1m block with the initial temperature
distribution is $T(x,0)=\sin\pi x$. The temperature on the boundaries are
constant. We set the boundary temperature to 0oC. The exact solution of this
problem is
$T(x,t)=e^{\pi^{2}\alpha t}\sin\pi x$
where $\alpha=k/\rho c_{p}=0.1$.
Figure 3: Heat conduction for 1D sinusoidal temperature
Figure 3 shows the SPH solution for heat conduction fitted with the exact
solution for times $t=0,0.1,0.2,...,1$. The result shows that the constant
temperature on the boundaries acts as heat sink of the system.
### 3.2 Three-dimensional liquid metal solidification
#### 3.2.1 Physical set up and material properties
The materials used in this simulation are liquid aluminium as the liquid metal
and steel as the mold. The geometry of the mold steel is cube-shaped with the
dimension 0.3m$\times$0.3m$\times$0.3m. The temperature of the mold steel is
200oC and the thermal conductivity is $30$ W/mK.
The mold is fully filled by the liquid aluminium. The liquid aluminium data is
taken from [12]. The solidus and liquidus temperatures are 536.1oC and 589.7oC
respectively. The viscosity of the liquid aluminium is temperature dependent.
The variation of the viscosity with temperature is shown in the Figure 4. The
viscosity above the liquidus temperature is 0.0118 Pa s and the viscosity
below the solidus temperature is 52.4 Pa s. The density of the liquid
aluminium is also temperature dependent. The liquid metal has a density of
2540 kg/m3 and the solid metal has a density of 2702 kg/m3. In the mushy zone,
the variation of density with temperature is close to linear. The variation of
the density is shown in the Figure 4. The other data for the aluminium
properties can be seen in Table LABEL:table:1.
Figure 4: Relation of material properties with temperature: (a) viscosity and (b) density Table 1: Thermal physical properties of aluminium Parameter | Value | Unit
---|---|---
Thermal conductivity of liquid | 94.14 | W/mK
Thermal conductivity of solid | 237.65 | W/mK
Specific heat of liquid | 1.084 | kJ/kg/K
Specific heat of solid | 0.963 | kJ/kg/K
Initial temperature | 600 | oC
Latent heat | 398 | kJ/kg
#### 3.2.2 Results and discussions
In this simulation, we use 15625 liquid metal particles with 0.0106154 m of
smoothing length. The simulation starts by arranging the liquid metal
particles position in the well-ordered position. Then the particles are
solidified rapidly through the side walls and the bottom wall. The results for
several time steps are shown in the Figure 5. Actually the domain is fully
filled with the liquid metal particles. However, we display the results from
the half part, so we can observe the inner part. The results show not only the
temperature distribution but also the phase transition. The temperature
distribution for the liquid phase and mushy phase are represented by the color
bar. The red color represents the liquid phase with temperature above 589.7oC
and the blue color represents the lowest temperature for the mushy phase that
is 536.1oC. The solid phase is represented by silver color where the
temperature is below 536.1oC.
Figure 5: Solidification process for several time step: (a) t = 0.0045 s, (b)
t = 0.5 s, (c) t = 1 s, (d) t = 1.5 s, (e) t = 2 s, and (f) t = 2.7272 s
The solidification process seems very fast. All of the liquid particles become
solid after 2.7272 s. It is because we set the initial temperature of a single
particle 600oC while the initial temperature of the boundary is 200oC. This
value is very close to the liquidus temperature that is 589.1oC. Another
reason is that the thermal conductivities of mushy and solid phase are very
high. So, if some particles are transformed to the solid phase, the heat
transfer will be faster. As an example, the temperature histories of a single
particle can be seen in Figure 6. It shows the different temperature drops at
different phases. After the particle reaches the mushy phase or solid phase,
the heat transfer works faster and leads the temperature to decrease rapidly.
Figure 6: Temperature history for a single particle. It reaches mushy phase at
$t=1.35455$ s and solid phase at $t=1.99545$ s
From the results in Figure 5, it can be seen that the volume of the liquid
metal is decreasing every time step and it leads to the shrinkage cavity
formation at the final solidification stage. It happened due to the difference
in density where the solid phase has a higher density than the liquid phase.
This affects the density of the mushy particles to continuously grow during
the cooling down process. Hence the equation of states in Eq.(8) gives
positive pressure; moreover the gravitational force also contributes to push
the particles to the bottom. This simulation seems to approximate nicely the
realistic phenomena. However, in the rapid cooling case, other defects such as
cracks or voids often appear. Presence of cracks means there are some empty
parts inside if we check the inner part of the particles distribution. The
particles distribution for different layers can be seen in Figure 7. We got
different patterns on different layers. There are some lower densities and
higher densities according to the number of particles. In the lower density
part some cracks should appear. However, in our cases the cracks do not
appear. It means that the model for liquid-solid interaction in 3D simulation
should be improved. Our model which assumes the solid group as fluid with high
viscosity and acts as undeformable object seems correct. It is because we can
produce the shrinkage cavity formation and the simulation is stable until all
the particle become solid. However, it is still not enough and needs to be
improved such as by applying some particular forces between mushy and solid
phases to produce other defects.
Figure 7: Density difference according to the number of particles at different
layers: (a) 3 $\leq$ z $\leq$ 6 cm, (b) 6 $\leq$ z $\leq$ 9 cm, and (c) 9
$\leq$ z $\leq$ 12 cm
## 4 Summary
In this paper, the behavior of liquid metal solidification has been studied
and simulated using smoothed particle hydrodynamics. The solidification
process was done by installing the enthalpy method in SPH formulation to solve
the heat transfer problem. The model of heat transfer process has been tested
to static particle grid in one dimension and it gives a relatively good
accuracy. The phase transition from liquid to solid was treated by considering
the effect of latent heat and nonisothermal phase change to the model of heat
transfer. The liquid-solid interaction model for 3D simulation also has been
performed. The results showed the appearance of a defect when liquid metal was
solidified. We got shrinkage cavity formation in the final solidification
stage. However, in the rapid cooling process, other defects often appear, to
obtain these defects, the model of liquid-solid interaction should be improved
such as by adding a particular force between mushy and solid phases. The
improvement of liquid-solid interaction model is left for the further study.
## References
* [1] H. Hu and S. A. Argyropoulos (1996). Mathematical modelling of solidification and melting: a review, Modelling Simul. Mater. Sci. Eng., 4, 4, pp. 371-396.
* [2] L. B. Lucy (1977), A numerical approach to the testing of the fission hypothesis, Astron. J., 82, 1013 –- 1024.
* [3] R. A. Gingold and J. J. Monaghan (1977), Smoothed particle hydrodynamics: theory and application to non–spherical stars, Mon. Not. R. Astr. Soc., 181, 375 – 389.
* [4] J. J. Monaghan, H. E. Huppert, and M. G. Worster (2005), Solidification using smoothed particle hydrodynamics, J. Comput. Phys., 206, 684 – 705.
* [5] P.W. Cleary, J. Ha, M. Prakash, and T. Nguyen (2010). Short shots and industrial case studies: Understanding fluid flow and solidification in high pressure die casting. Applied Mathematical Modelling. 34, 2018 -– 2033.
* [6] M. Liu and G. Liu (2010). Smoothed Particle Hydrodynamics (SPH): an overview and recent developments, Archives of Computational Methods in Engineering, 17, 1, pp. 25 – 76.
* [7] J. F. Wendt (2009), Computational fluid dynamics: An introduction. Springer, Berlin.
* [8] J.P. Morris, J.J. Monaghan (1997). A switch to reduce SPH viscosity, J. Compt. Phys., 136, 41 –- 50.
* [9] P.W. Cleary (1998). Modelling confined multi-material heat and mass flows using SPH. Applied Mathematical Modelling. 22, 981 -– 993.
* [10] P.W. Cleary and J.J. Monaghan (1999). Conduction modelling using smoothed particle hydrodynamics, J. Comput. Phys., 148, 1, pp. 227 – 264.
* [11] B. Nedjar (2001) An enthalpy-based finite element method for nonlinear heat problems involving phase change, Computers and structures, 80, 9 – 21.
* [12] P.W Cleary (2010). Extension of SPH to predict feeding, freezing and defect creation in low pressure die casting. Applied Mathematical Modelling. 34, 3189 -– 3201.
|
arxiv-papers
| 2013-09-17T09:07:50 |
2024-09-04T02:49:50.993943
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Raden Ahnaf Faqih S, Christian Fredy Naa",
"submitter": "Christian Naa",
"url": "https://arxiv.org/abs/1309.4234"
}
|
1309.4251
|
# Optimal Distributed Controller Design with
Communication Delays: Application to Vehicle Formations
Hamid Reza Feyzmahdavian, Assad Alam, and Ather Gattami H. R. Feyzmahdavian is
with Electrical Engineering, ACCESS Linnaeus Centre, Royal Institute of
Technology, 100 44 Stockholm, Sweden. [email protected]. Alam is with Research
and Development, Scania CV AB, 151 87 Södertälje, Sweden.
[email protected]. Gattami is with Electrical Engineering, ACCESS
Linnaeus Centre, Royal Institute of Technology, 100 44 Stockholm, Sweden.
[email protected] work was supported by Scania CV AB, VINNOVA - FFI, and the
Swedish Research Council.
###### Abstract
This paper develops a controller synthesis algorithm for distributed LQG
control problems under output feedback. We consider a system consisting of
three interconnected linear subsystems with a delayed information sharing
structure. While the state-feedback case of this problem has previously been
solved, the extension to output-feedback is nontrivial, as the classical
separation principle fails. To find the optimal solution, the controller is
decomposed into two independent components. One is delayed centralized LQR,
and the other is the sum of correction terms based on additional local
information. Explicit discrete-time equations are derived whose solutions are
the gains of the optimal controller.111A preliminary version of this work was
presented in [1].
## I INTRODUCTION
The systems to be controlled are in many application domains getting larger
and more complex. When there is interconnection between different dynamical
systems, conventional optimal control algorithms provide a solution where
centralized state information is required. However, it is often preferable and
sometimes necessary to have a distributed control structure, since in many
practical problems, the physical or communication constraints often impose a
specific interconnection structure. Hence, it is interesting to design
distributed feedback controls for systems of a certain structure and examine
their overall performance.
The control problem and methodology in this paper is motivated by systems
involving a chain of closely spaced heavy duty vehicles (HDVs), generally
referred to as vehicle platooning. The objective is to maintain a predefined
headway to the vehicle ahead, while maintaining safety and minimizing the fuel
consumption. Information technology is paving its path into the transport
industry, enabling the possibility of automated control strategies. Governing
vehicle platoons by an automated control strategy, the overall traffic flow is
expected to improve [2] and the road capacity will increase significantly [3],
without endangering safety [4]. By traveling at a close intermediate spacing
the air drag is reduced for each vehicle in the platoon. Thereby, the control
effort and inherently the fuel consumption can be reduced significantly [5].
This creates a coupling of the dynamics between neighboring vehicles
throughout the platoon. However, as the intermediate spacing is reduced the
control becomes tighter due to safety aspects; mandating an increase in
control action and inherently the fuel consumption through additional
acceleration and braking. The fuel consumption constitutes approximately 30 %
of the overall cost for a fleet owner [6]. Hence, it is of vast interest for
the industry to find a fuel optimal control. Considering the physical
constraints in radio, it cannot be assumed that state information is available
at every instance in time. Thus, a distributed control strategy is crucial for
practical implementation.
In recent work [7, 8, 9], distributed control has been studied under the
assumption of spatial invariance. Control for chain structures in the context
of platoons has been studied through various perspectives, e.g., [10, 11, 12,
13, 14]. It has been shown that control strategies may vary depending on the
available information within the platoon. Maintaining a suitable relative
distance, stability and robustness of the platoon have been identified to be
amongst the main criteria to be considered. However, communication constraints
have not in general been considered in control design for platooning
applications and the controllers have mainly been ad hoc by tuning the control
parameters. In [15, 16], linear quadratic Gaussian (LQG) control under
appropriate assumptions on communication delays between the controllers was
considered. While a computationally efficient solution was presented for a
sequence of vehicles moving in formation, the controller structure is not
provided by the corresponding semi-definite programming. A structured
sequential design was introduced in [17], where the preceding vehicle’s
dynamics along with its states were conveyed through wireless communication.
It resulted in a suboptimal control strategy, where physical coupling to a
follower vehicle and communication delays were not considered. Mounted radar
sensors allows each vehicle to measure the relative distance and velocity of
the preceding vehicle. Additional information, providing local information,
has lately been introduced through wireless information. However, wireless
systems introduce information delays to the system in certain cases due to
limitations in radio. Furthermore, varying external environment factors impose
process disturbances on the system.
In this work, we are primarily concerned with forming a distributed control,
that accounts for the interconnection between neighboring vehicles, correlated
process disturbances, as well as communication delays. The control is solely
based on local model knowledge, over the class of LQG control for chain
structured interconnection graphs. The received information is assumed to be
common after two time step delays.
The main contribution of this paper is to derive an LQG controller, which is
easy to implement and optimal under a delayed information sharing pattern for
chain structures. In addition to communication delays, the distributed optimal
control is based upon systems with interconnected dynamics to both neighboring
vehicles and local state information. Derived from the characteristics of
actual Scania HDV’s, we present a discrete system model that includes physical
coupling with both neighboring vehicles. We also investigate the performance
of the proposed controllers, under normal operating conditions for an HDV
platoon, with respect to physical constraints that are imposed in a practical
set-up.
The outline of the paper is as follows. The general system and problem
description is given in Sec. II, which in turn determines the structure of the
optimal controller. The theoretical premise for the optimal controller is
presented in Sec. IV, where it is shown that the problem can be decomposed
into two separate optimization problems. Finally, we evaluate the performance
of the derived controller through numerical results in Sec. V and give
concluding remarks in Sec. VI.
### Notation
Throughout the paper, we use the following notation: matrices are written in
uppercase letters and vectors in lowercase letters. The $i^{th}$ component of
a vector $x$ is denoted by $x_{i}$. Let $[x]_{S}$ be the sub-vector of $x$
containing only those components with indices in set $S$. For instance, if
$S=\\{1,3\\}$, then $[x]_{S}$ is given by
$[x]_{S}=\begin{bmatrix}x_{1}&x_{3}\end{bmatrix}^{T}$. The sequence $x(0)$,
$x(1)$, $\ldots$ , $x(k)$ is denoted by $x(0:k)$.
diag$(x)$ denotes a diagonal matrix whose diagonal elements are given by those
of the enclosed vector $x$. Let $X$ be a matrix partitioned into blocks. We
use $[X]_{ij}$ and $[X]_{i}$ to represent the block in block position $ij$ and
$i$th block row, respectively. $[X]_{S_{1}S_{2}}$ denotes the sub-matrix of
$X$ containing exactly those rows and columns corresponding to the sets $S1$
and $S2$, respectively. For instance
$[X]_{\\{1\\}\\{2,3\\}}=\begin{bmatrix}X_{12}&X_{13}\end{bmatrix}$. The trace
of square matrix is denoted by $\textbf{Tr}\\{X\\}$. We use $X^{+}$ and
$X^{-}$ to represent $X(k+1)$ and $X(k-1)$ respectively, when appropriate.
Given $A\in\mathbb{R}^{m\times n}$, we can write $A$ in terms of its columns
as $A=\begin{bmatrix}a_{1}&\cdots&a_{n}\end{bmatrix}$. The operation vec$(A)$
results in a $mn\times 1$ column vector
$\textrm{vec}(A)=\begin{bmatrix}a_{1}^{T}&\cdots&a_{n}^{T}\end{bmatrix}^{T}$.
We denote by $\textrm{vec}^{\star}(A)$, the sub-vector of $\textrm{vec}(A)$
containing only nonzero elements. Let $A\in\mathbb{R}^{m\times n}$ and
$B\in\mathbb{R}^{r\times s}$, then the operation $A\otimes
B\in\mathbb{R}^{mr\times ns}$ denotes the Kronecker product of $A$ and $B$.
We denote the expectation of a random variable $x$ by $\textbf{E}\\{x\\}$. The
conditional expectation of $x$ given $y$ is denoted by $\textbf{E}\\{x|y\\}$.
## II System Model and Problem Description
In this section we present the physical properties of the system that we are
considering. We state the nonlinear dynamics of a single vehicle and the model
for the aerodynamics, which induces the physical coupling. Then we present the
linear discrete system model for a heterogeneous HDV platoon and its
associated cost function. The communication constraints and physical coupling
is then used to motivate the structure of the controller. Finally, the problem
formulation is given.
### II-A System Model
We consider an HDV platoon as depicted in Fig. 1. The state equation of a
single HDV is modeled as,
$\displaystyle\dot{s}$ $\displaystyle=v,$ (1) $\displaystyle m_{t}\dot{v}$
$\displaystyle=F_{engine}-F_{brake}-F_{airdrag}(v)$ $\displaystyle\hskip
44.3863pt-F_{roll}(\alpha)-F_{gravity}(\alpha),$
$\displaystyle=k_{u}u-k_{b}F_{brake}-k_{d}v^{2}$ $\displaystyle\hskip
33.0052pt-k_{fr}\cos\alpha-k_{g}\sin\alpha,$
where $v$ is the vehicle velocity, $m_{t}$ denotes the accelerated mass and
$u\in\mathbb{R}$ denotes the net engine torque. $k_{u},k_{b},k_{d},k_{fr}$,
and $k_{g}$ denote the characteristic vehicle and environment coefficients for
the engine, brake, air drag, road friction, and gravitation respectively.
Figure 1: The figure shows a platoon of $M$ heavy duty vehicles, where each
vehicle is able to communicate with its neighbors.
The variation in aerodynamics between the vehicles is essential in the
analysis of fuel reduction potential for HDVs. For a single HDV it can amount
up to 50 % of the total resistive forces at full speed. It is significantly
reduced when operating in a platoon formation and a coupling between the
vehicles is induced. To account for the aerodynamics, the air drag
characteristic coefficient in (1) can be modeled as [18]
$\displaystyle\begin{split}\tilde{k}_{d}&=k_{d}(1-\frac{\Phi(d)}{100}-\frac{\phi(d)}{100}),\\\
\Phi(d)&=\alpha_{1}d+\alpha_{2},0\leq d\leq 60\\\
\phi(d)&=\beta_{1}d+\beta_{2},0\leq d\leq 15\end{split}$
where $d$ is the longitudinal relative distance between two vehicles,
$\Phi(d)$ and $\phi(d)$ are linear piecewise affine functions of the change in
air drag due to a preceding and a follower vehicle respectively, and
$\alpha_{1},\alpha_{2},\beta_{1},\beta_{2}$ are positive constants. The
relative distance reference could be constant or, as in this case, time
varying. It is determined by setting a desired time gap $\tau$ s, which in
turn determines the spacing policy as $d_{ref}(k)={\tau}v(k)$. Thereby, the
vehicles will maintain a larger intermediate spacing at higher velocities.
When studying the behavior of an HDV platoon, the velocity does not deviate
significantly from the lead vehicle’s velocity. Ideally, all vehicles should
maintain a constant speed and intermediate distance. Thus, a linearized model
should give a sufficient description of the system behavior under these
conditions. By linearizing and applying a one step forward discretization to
(1), the discrete model for an HDV platoon with respect to a set reference
velocity, an engine torque which maintains the velocity, a fixed spacing
between the vehicles, and a constant slope is hence given by
$x(k+1)=Ax(k)+Bu(k)+w(k),$
where
$\displaystyle A$
$\displaystyle=\begin{bmatrix}{\Theta_{1}}&\gamma_{2}&0&0&0&\cdots&0&0&0\\\
1&1&-1&0&0&\cdots&0&0&0\\\
0&\delta_{2}&\Theta_{2}&\gamma_{3}&0&\cdots&0&0&0\\\
0&0&1&1&-1&\cdots&0&0&0\\\
0&0&0&\delta_{3}&\Theta_{3}&\cdots&\gamma_{4}&0&0\\\
\vdots&\vdots&\vdots&\vdots&\vdots&\ddots&\vdots&\vdots&\vdots\\\
0&0&0&0&0&\cdots&\Theta_{M-1}&\gamma_{M-1}&0\\\ 0&0&0&0&0&\cdots&1&1&-1\\\
0&0&0&0&0&\cdots&0&\delta_{M}&\Theta_{M}\\\ \end{bmatrix},$
$\displaystyle\begin{split}B&=\begin{bmatrix}k_{u_{1}}&0&0&\cdots&0\\\
0&0&0&\cdots&0\\\ 0&k_{u_{2}}&0&\cdots&0\\\ 0&0&0&\cdots&0\\\
0&0&k_{u_{3}}&\cdots&0\\\ \vdots&\vdots&\vdots&\ddots&\vdots\\\
0&0&0&\cdots&0\\\ 0&0&0&\cdots&k_{u_{M}}\end{bmatrix},\quad
x=\begin{bmatrix}{v_{1}}\\\ {d_{12}}\\\ {v_{2}}\\\ {d_{23}}\\\ {v_{3}}\\\
\vdots\\\ {v_{M-1}}\\\ {d_{(M-1)M}}\\\ {v_{M}}\end{bmatrix},\\\
u&=\begin{bmatrix}{u_{1}}\\\ {u_{2}}\\\ {u_{3}}\\\ \vdots\\\
{u_{M}}\end{bmatrix},\begin{array}[]{rl}\Theta_{1}&=1-T_{s}2k_{d}(d_{0})v_{0}/m_{t},\\\
\Theta_{i}&=1-T_{s}2k_{d}\Phi{(d_{0})}v_{0}/m_{t},\quad i=2,\dots,M,\\\
\delta_{i}&=-T_{s}\alpha_{1}k_{d}v_{0}^{2}/m_{t},\\\
\gamma_{i}&=-T_{s}\beta_{1}k_{d}v_{0}^{2}/m_{t},\end{array}\end{split}$ (2)
and $T_{s}$ is the sampling time. Thus, the system has a block diagonal
structure and can be grouped into subsystems as indicated in (2). The general
representation of the derived system can be stated as
$\displaystyle\begin{bmatrix}x_{1}(k+1)\\\ x_{2}(k+1)\\\ x_{3}(k+1)\\\
\vdots\\\ x_{M}(k+1)\end{bmatrix}=$
$\displaystyle\begin{bmatrix}A_{11}&A_{12}&0&\cdots&0\\\
A_{21}&A_{22}&A_{23}&\cdots&0\\\ 0&A_{32}&A_{33}&\cdots&0\\\
\vdots&\vdots&\vdots&\ddots&\vdots\\\
0&0&0&\cdots&A_{MM}\end{bmatrix}\begin{bmatrix}x_{1}(k)\\\ x_{2}(k)\\\
x_{3}(k)\\\ \vdots\\\ x_{M}(k)\end{bmatrix}$ $\displaystyle+$
$\displaystyle\begin{bmatrix}B_{1}&0&0&\cdots&0\\\ 0&B_{2}&0&\cdots&0\\\
0&0&B_{3}&\cdots&0\\\ \vdots&\vdots&\vdots&\ddots&\vdots\\\
0&0&0&\cdots&B_{M}\\\ \end{bmatrix}\begin{bmatrix}u_{1}(k)\\\ u_{2}(k)\\\
u_{3}(k)\\\ \vdots\\\ u_{M}(k)\end{bmatrix}+w(k)$ (3)
where the corresponding vehicle states for each subsystem are
$x_{1}(k)=v_{1}(k),\quad x_{i}(k)=\begin{bmatrix}d_{i-1,i}\\\
v_{i}\end{bmatrix},\quad i=2,\dots,M.$
In practice, many random disturbances are imposed upon a vehicle in motion.
The varying road topology has a strong impact due the extensive mass of the
HDVs. Weather conditions might vary and traffic conditions might change.
Furthermore, variation in wind affects all the vehicles in the platoon and
therefore the process noise is considered to be correlated. Hence, the
disturbance, $w(k)$ in (3), is assumed to be a Gaussian white noise with a
full positive definite covariance matrix $W$. We also assume that the initial
state $x(0)$ is uncorrelated with $w(k)$ for all $k$, with zero mean and
covariance matrix $P_{0}$.
While a general problem was defined, for simplicity, consider an $M=3$ HDV
platoon. In this case, the dynamics of the system given in (4) is
$\displaystyle x_{1}(k+1)$
$\displaystyle=A_{11}x_{1}(k)+A_{12}x_{2}(k)+B_{1}u_{1}(k)$ $\displaystyle
x_{2}(k+1)$
$\displaystyle=A_{21}x_{1}(k)+A_{22}x_{2}(k)+A_{23}x_{3}(k)+B_{2}u_{2}(k)$
$\displaystyle x_{3}(k+1)$
$\displaystyle=A_{32}x_{2}(k)+A_{33}x_{3}(k)+B_{3}u_{3}(k)$ (4)
It can be seen in (4) that the state of vehicle $1$ is affected by the states
of vehicle $2$ in the next time step. Whereas, the state of vehicle $1$
affects the states of $3$ after two time steps, through vehicle $2$. Vehicle
$2$ on the other hand is affected by both vehicle $1$ and $3$ in the next time
step.
The local models can be conveyed at a single point in time between each
subsystem, through wireless communication. However, the system is time
critical due to safety aspects and communication should be kept at minimum so
the channel is not congested and latency is introduced. Assume that passing
information from one vehicle to another vehicle takes one time step, so the
available information set of each vehicle at time $k$ can be described as
$\displaystyle\mathcal{I}_{1}(k)$
$\displaystyle=\\{x_{1}(k),x_{1}(k-1),x_{2}(k-1),x(0:k-2)\\}$
$\displaystyle\mathcal{I}_{2}(k)$
$\displaystyle=\\{x_{2}(k),x(k-1),x(0:k-2)\\}$
$\displaystyle\mathcal{I}_{3}(k)$
$\displaystyle=\\{x_{3}(k),x_{2}(k-1),x_{3}(k-1),x(0:k-2)\\}$ (5)
The three vehicles share all past information with two-step communication
delay, as described in (5). The assumptions about the information structure
and the sparsity of dynamics guarantee that information propagates at least as
fast as the dynamics. This information pattern is a simple case of partially
nested information structure. It is shown in [19] that if the information
structure is partially nested, then the optimal controller exists, it is
unique, and linear. Therefore, the optimal controller for three vehicles under
the given information set has the form
$\displaystyle u_{1}(k)$ $\displaystyle=$ $\displaystyle
f_{11}\bigl{(}x_{1}(k)\bigr{)}+f_{12}\bigl{(}x_{1}(k-1),x_{2}(k-1)\bigr{)}$
$\displaystyle\hskip 49.22314pt+f_{13}\bigl{(}x(0:k-2)\bigr{)}$ $\displaystyle
u_{2}(k)$ $\displaystyle=$ $\displaystyle
f_{21}\bigl{(}x_{2}(k)\bigr{)}+f_{22}\bigl{(}x(k-1)\bigr{)}$
$\displaystyle\hskip 49.22314pt+f_{23}\bigl{(}x(0:k-2)\bigr{)}$ $\displaystyle
u_{3}(k)$ $\displaystyle=$ $\displaystyle
f_{31}\bigl{(}x_{3}(k)\bigr{)}+f_{32}\bigl{(}x_{2}(k-1),x_{3}(k-1)\bigr{)}$
(6) $\displaystyle\hskip 49.22314pt+f_{33}\bigl{(}x(0:k-2)\bigr{)}$
where $f_{ij}$ denotes a linear function in all its variables. Consequently,
the optimal control $u(k)$ can be expressed as
$\displaystyle u(k)=F(k)x(k)+G(k){x(k-1)}+f\bigl{(}{x}(0:k-2)\bigr{)}$ (7)
where $f={\begin{bmatrix}f_{13}^{T}&f_{23}^{T}&f_{33}^{T}\\\
\end{bmatrix}}^{T}$ and
$\displaystyle F(k)$ $\displaystyle=\begin{bmatrix}F_{11}&0&0\\\ 0&F_{22}&0\\\
0&0&F_{33}\\\ \end{bmatrix},\;G(k)$
$\displaystyle=\begin{bmatrix}G_{11}&G_{12}&0\\\ G_{21}&G_{22}&G_{23}\\\
0&G_{32}&G_{33}\\\ \end{bmatrix}.$
### II-B Cost Function
The objective of the lead vehicle is to minimize the fuel consumption and
control input, while maintaining a set reference velocity. The objective of
the follower vehicles in addition is to follow the preceding vehicles
velocity, while maintaining a set intermediate spacing. Hence, similar to what
we presented for the continuous LQR in [17], the weights for a $M$ HDV platoon
can be set up based upon the performance objectives as
$\displaystyle J(u^{*})=$
$\displaystyle\min_{u}~{}\sum_{k=0}^{N-1}\Big{(}\sum_{i=2}^{M}w_{i}^{\tau}(d_{(i-1)i}(k)-\tau
v_{i}(k))^{2}$ $\displaystyle\hskip
38.41121pt+w_{i}^{\Delta{v}}(v_{i-1}(k)-v_{i}(k))^{2}$ $\displaystyle\hskip
38.41121pt+w_{i}^{d}d_{(i-1)i}^{2}(k)+\sum_{i=1}^{M}w_{i}^{v}v_{i}^{2}(k)+w_{i}^{u_{i}}u_{i}^{2}(k)\Big{)}$
$\displaystyle=$
$\displaystyle\min_{u}~{}\sum_{k=0}^{N-1}\sum_{i=2}^{M}\begin{bmatrix}v_{i-1}(k)\\\
d_{(i-1)i}(k)\\\ v_{i}(k)\end{bmatrix}^{T}Q_{i}\begin{bmatrix}v_{i-1}(k)\\\
d_{(i-1)i}(k)\\\ v_{i}(k)\end{bmatrix}+R_{i}u_{i}^{2}(k)$ $\displaystyle\hskip
38.41121pt+w^{v_{1}}v_{1}^{2}(k)+w^{u_{1}}u_{1}^{2}(k)$ (8)
where
$\displaystyle\begin{split}&Q_{i}=\begin{bmatrix}w_{i}^{\Delta{v}}&0&-w_{i}^{\Delta{v}}\\\
0&w_{i}^{d}+w_{i}^{\tau}&-\tau{w}_{i}^{\tau}\\\
-w_{i}^{\Delta{v}}&-\tau{w}_{i}^{\tau}&\tau^{2}w_{i}^{\tau}+w_{i}^{\Delta{v}}+w_{i}^{v}\end{bmatrix},\\\
&Q_{1}=\begin{bmatrix}w^{v_{1}}&0\\\
0&w^{u_{1}}\end{bmatrix},R_{i}=w_{i}^{u_{i}}.\end{split}$
The weights in (8) give a direct interpretation of how to enforce the
objectives for a vehicle traveling in a platoon. The value of $w_{i}^{\tau}$
determines the importance of not deviating from the desired time gap. Hence, a
large $w_{i}^{\tau}$ puts emphasis on safety. $w_{i}^{\Delta{v}}$ creates a
cost for deviating from the velocity of the preceding vehicle, and
$w_{i}^{u_{i}}$ punishes the control effort which is proportional to the fuel
consumption. The following terms, $w_{i}^{d},w_{i}^{v}$, put a cost on the
deviation from the linearized states. Note that the main objective is to
maintain a set intermediate distance, while maintaining a fuel efficient
behavior. Therefore, $w_{i}^{\tau},w_{i}^{\Delta{v}}$ and $w_{i}^{u_{i}}$ must
be set larger than the remaining weights.
### II-C Problem Formulation
We consider a HDV platooning scenario where each vehicle only receives
information regarding the relative position and velocity of the immediate
neighboring vehicles. The objective is to design a controller that can handle
a two time step delay.
The aim is to utilize the given structure of the considered system, where we
want to minimize the cost function
$\displaystyle\begin{split}J=&\textbf{E}\\{x(N)^{T}Q_{0}x(N)\\}\\\
&+\sum_{k=0}^{N-1}\textbf{E}\\{x^{T}(k)Qx(k)+u^{T}(k)Ru(k)\\},\end{split}$ (9)
subject to the sparse system dynamics in (3) and the performance objectives in
(8). The primary difficulty arises from the imposed information constraints
given in (5).
Thus, the problem that we solve in this paper is finding an analytical
expression for an optimal control input $u_{i}(k)$, which must be a function
of the admissible information set $\mathcal{I}_{i}(k)$, where each subsystem
control input is unique and a linear function denoted as
$\displaystyle
u_{i}(k)=\mu_{i}\bigl{(}\mathcal{I}_{i}(k)\bigr{)},\;i=1,\dots,M.$ (10)
###### Assumption 1
The matrices $Q_{0}$ and $Q$ in (9) are positive semi-definite, and $R$ is
positive definite.
## III Main Result
In this section we present the optimal controller for three-vehicle problem.
The proof for this result is presented in the remaining sections.
###### Theorem 1
_Suppose that $W$ is positive definite and that Assumption $1$ holds. Define
the matrix $D\triangleq\begin{bmatrix}F&M\end{bmatrix}$ where $M$ has the same
sparsity structure as $G$. Let $S$ be the index set of non-zero elements of_
$\text{vec}(D)$,
$S\triangleq\left\\{i:\;\textrm{vec}_{i}(D)\neq 0\right\\}.$
_Suppose there exists a stabilizing solution $X$ to the algebraic Riccati
equation_
$\displaystyle X=A^{T}XA+Q+A^{T}XB(B^{T}XB+R)^{-1}B^{T}XA$
_We then define_
$\displaystyle H$ $\displaystyle=B^{T}XB+R\vspace{0.1cm}$ $\displaystyle L$
$\displaystyle=(B^{T}XB+R)^{-1}B^{T}XA\vspace{0.1cm}$
_and let_
$\displaystyle Y$
$\displaystyle=\begin{bmatrix}W\otimes(H+B^{T}L^{T}HLB)&-W\otimes
B^{T}L^{T}H\\\ -W\otimes HLB&W\otimes H\end{bmatrix}$ $\displaystyle b$
$\displaystyle=\begin{bmatrix}W\otimes H\\\
0\end{bmatrix}\mbox{vec}(L)+\begin{bmatrix}-W\otimes B^{T}L^{T}H\\\ W\otimes
H\end{bmatrix}\textrm{vec}(LA)$
_Then, the optimal controller gains are given by:_
$\displaystyle\textrm{vec}^{\star}(F)$
$\displaystyle=\begin{bmatrix}I&0\end{bmatrix}[Y]_{SS}^{-1}[b]_{S}\vspace{1mm}$
$\displaystyle\textrm{vec}^{\star}(M)$
$\displaystyle=\begin{bmatrix}0&I\end{bmatrix}[Y]_{SS}^{-1}[b]_{S}$
_and the optimal controller has the realization_
$\displaystyle\zeta(k+1)=$ $\displaystyle Ax(k)+Bu(k)$
$\displaystyle\xi(k+1)=$ $\displaystyle
A\zeta(k)+BM(x(k-1)-\zeta(k-1))+BL\xi(k)$ $\displaystyle u(k)=$ $\displaystyle
F(x(k)-\zeta(k))$ $\displaystyle+M(x(k-1)-\zeta(k-1))+L\xi(k)$
Note that blocks of matrices $F$ and $M$ can be computed from the
$\textrm{vec}^{\star}(F)$ and $\textrm{vec}^{\star}(G)$, respectively. For
example,
$\textrm{vec}^{\star}(F)=\textrm{vec}\left(\begin{bmatrix}F_{11}&F_{22}&F_{33}\end{bmatrix}\right)$.
It will be shown that $\xi(k)$ is the minimum-mean square estimate of $x(k)$
given the common information $x(0:k-2)$; that is,
$\xi(k)=\textbf{E}\\{x(k)|x(0:k-2)\\}$. Thus, the optimal controller of three-
vehicle problem is the centralized $LQR$ controller under the classical
information structure with two-step delay plus correction terms based on the
local information at time $k$.
## IV Optimal Controller Derivation
In this section, we present the preliminary lemmas that are used to prove the
results in Theorem $1$. Before proceeding further, we need to state the
following proposition which later permits us to decompose $J$ into two
separate parts.
###### Proposition 1 ([20])
_Define the matrices_
$\displaystyle X(k)=$ $\displaystyle A^{T}X^{+}A+Q$ (11)
$\displaystyle-(A^{T}X^{+}B)(B^{T}X^{+}B+R)^{-1}(B^{T}X^{+}A)$ $\displaystyle
H(k)=$ $\displaystyle B^{T}X^{+}B+R\vspace{0.1cm}$ $\displaystyle L(k)=$
$\displaystyle(B^{T}X^{+}B+R)^{-1}B^{T}X^{+}A\vspace{0.1cm}$
_for $k=0,\cdots,N-1$ and where $X(N)=Q_{0}$. Then the cost function (9) can
be written as _
$\displaystyle J=$
$\displaystyle\underbrace{\sum_{k=0}^{N-1}\textbf{E}\left\\{\bigl{(}u(k)-L(k)x(k)\bigr{)}^{T}H(k)\bigl{(}u(k)-L(k)x(k)\bigr{)}\right\\}}_{J_{u}}$
$\displaystyle+\underbrace{x^{T}(0)X(0)x(0)+\sum_{k=0}^{N-1}\textbf{Tr}\\{X(k+1)W\\}}_{J_{w}}$
_where both the zero-mean property of $w(k)$ and independence of $w(k)$ and
$(x(k),u(k))$ are exploited. Moreover, $J_{w}$ is independent of $u$._
From Proposition 1, it can be seen that minimizing $J$ is equivalent to
minimizing $J_{u}$. Note that, under the Assumption 1, $H(k)$ is positive
definite.
### IV-A State Decomposition
The first step towards finding the optimal controller is decomposing the state
vector into independent terms.
###### Lemma 1
_The state vector can be decomposed as_
$\displaystyle x(k)=$
$\displaystyle\underbrace{w(k-1)+\bigl{(}A+BF(k-1)\bigr{)}w(k-2)}_{x^{1}(k)}$
$\displaystyle+\underbrace{\textbf{E}\left\\{x(k)|x(0:k-2)\right\\}}_{x^{2}(k)}$
_where $x^{1}(k)$ and $x^{2}(k)$ are independent random variables._
###### Proof:
The term $x^{2}(k)$ is the conditional estimate of state $x(k)$ given the
piece of information shared between all vehicles, and $x^{1}(k)$ is the
estimation error. The independence between $x(k)-x^{2}(k)$ and $x^{2}(k)$ can
be established by Proposition $4$b in the appendix. To calculate $x^{1}(k)$,
we proceed in three steps. First consider
$x(k-1)=Ax(k-2)+Bu(k-2)+w(k-2)$
Since $x(k-2)$ belongs to the sequence $x(0:k-2)$ and $u(k-2)$ is a
deterministic function of $x(0:k-2)$, we have
$x(k-1)-\textbf{E}\\{x(k-1)|x(0:k-2)\\}=w(k-2)$ (12)
where we used the zero-mean and independence of $w(k-2)$ and $x(0:k-2)$. The
structure of controller is given by equation (7), so $u(k-1)$ can be written
as
$u(k-1)=F(k-1)x(k-1)+G(k-1){x(k-2)}+f\bigl{(}{x}(0:k-3)\bigr{)}$
Since $G(k-1){x(k-2)}+f\bigl{(}{x}(0:k-3)$ is a deterministic function of
$x(0:k-2)$, we have
$\displaystyle u(k-1)-\textbf{E}\\{u(k-1)|x(0:k-2)\\}$ (13) $\displaystyle=$
$\displaystyle F(k-1)\bigl{(}x(k-1)-\textbf{E}\\{x(k-1)|x(0:k-2)\\}\bigr{)}$
$\displaystyle=$ $\displaystyle F(k-1)w(k-2)$
where we substituted (12) into the second line. Finally, note that $w(k-1)$
and $x(0:k-2)$ are independent. Therefore,
$\displaystyle x(k)$ $\displaystyle-\textbf{E}\\{x(k)|x(0:k-2)\\}$
$\displaystyle=$ $\displaystyle
w(k-1)+A\left(x(k-1)-\textbf{E}\\{{x}(k-1)|x(0:k-2)\\}\right)$
$\displaystyle\hskip
38.41139pt+B\left(u(k-1)-\textbf{E}\\{{u}(k-1)|x(0:k-2)\\}\right)$
$\displaystyle=$ $\displaystyle w(k-1)+\bigl{(}A+BF(k-1)\bigr{)}w(k-2)$ (14)
where we substituted (12) into the second line and (13) into the third line.
Thus the result follows. ∎
### IV-B Controller Decomposition
Now that the state has been decomposed into two independent terms, the control
input $u(k)$ can be decomposed in a similar fashion.
###### Lemma 2
_The control input $u(k)$ can be decomposed as_
$\displaystyle u(k)$
$\displaystyle=\underbrace{F(k)w(k-1)+M(k)w(k-2)}_{u^{1}(k)}+u^{2}(k)$
_where $u^{1}(k)$ and $u^{2}(k)$ are independent, $u^{2}(k)$ is a linear
function of $x(0:k-2)$, and_
$M(k)=F(k)\left(A+BF(k-1)\right)+G(k).$
###### Proof:
Let $u^{2}(k)=\textbf{E}\\{u(k)|x(0:k-2)\\}$, then $u^{2}(k)$ is a linear
function of $x(0:k-2)$ and independent of $u(k)-u^{2}(k)$. Note that
$f(x(0:k-2))$ is a linear function of $x(0:k-2)$, so $u^{1}(k)$ is computed as
$\displaystyle u^{1}(k)=$ $\displaystyle u(k)-\textbf{E}\\{u(k)|x(0:k-2)\\}$
$\displaystyle=$ $\displaystyle
F(k)\bigl{(}x(k)-\textbf{E}\\{x(k)|x(0:k-2)\\}\bigr{)}$
$\displaystyle+G(k)\bigl{(}x(k-1)-\textbf{E}\\{x(k-1)|x(0:k-2)\\}\bigr{)}$
$\displaystyle=$ $\displaystyle F(k)(w(k-1)+\left(A+BF(k-1)\right)w(k-2))$
$\displaystyle+G(k)w(k-2)$
where we used equation (7) in the first line, (14) in the second line and (12)
in the third line. The proof is completed by defining
$M(k)=F(k)\left(A+BF(k-1)\right)+G(k)$. ∎
###### Remark 1
Since $B$ and $F$ are diagonal matrices, $G(k)$ and $F(k)A$ have the same
sparsity structures. Therefore, sparsity structure of $M(k)$ and $G(k)$ are
also the same.
From Lemmas 1 and 2, $x^{2}(k)$ and $u^{2}(k)$ are linear functions of
$x(0:k-2)$ which is independent of $x^{1}(k)$ and $u^{1}(k)$. As a result the
cost function $J_{u}$ can be decomposed as:
$\displaystyle J_{u}$
$\displaystyle=\underbrace{\sum_{k=0}^{N-1}\textbf{E}\left\\{\bigl{(}u^{1}(k)-L(k)x^{1}(k)\bigr{)}^{T}H(k)\bigl{(}u^{1}(k)-L(k)x^{1}(k)\bigr{)}\right\\}}_{J^{1}_{u}}$
$\displaystyle+\underbrace{\sum_{k=0}^{N-1}\textbf{E}\left\\{\bigl{(}u^{2}(k)-L(k)x^{2}(k)\bigr{)}^{T}H(k)\bigl{(}u^{2}(k)-L(k)x^{2}(k)\bigr{)}\right\\}}_{J^{2}_{u}}$
The advantage of this decomposition of $J_{u}$ is that we now have two
subproblems on the form
$\displaystyle\begin{split}\min_{u^{1}}~{}&J^{1}_{u}(x^{1},u^{1})\\\
\mbox{subject to}~{}&u^{1}(k)=F(k)w(k-1)+M(k)w(k-2)\end{split}$ (15)
$\displaystyle\min_{u^{2}}~{}J^{2}_{u}(x^{2},u^{2})$
$\displaystyle\mbox{subject to}~{}u^{2}(k)=f\bigl{(}x(0:k-2)\bigr{)}$ (16)
### IV-C Finite Horizon Controller Derivation
First consider minimization problem (16). Before proceeding, let us state the
following proposition which allows us to find the optimal control $u^{2}(k)$.
###### Proposition 2 ([20])
_Consider the discrete time linear system_
$x(k+1)=Ax(k)+Bu(k)+w(k)$
_where $w(k)$ is a zero mean Gaussian white noise. Assume that
$u(k)=\mu\bigl{(}x(0:k)\bigr{)}$. Then the optimal control which minimizes the
cost function $J_{u}$, is given by_
$u(k)=L(k)x(k)$
The mapping from $x^{2}(k)$ to $u^{2}(k)$ is given in the following lemma.
###### Lemma 3
_The dynamics of $x^{2}$ can be written as_
$x^{2}(k+1)=Ax^{2}(k)+Bu^{2}(k)+T(k)w(k-2)$ (17)
_where $T(k)=A(A+BF(k-1))+BM(k)$._
###### Proof:
See appendix. ∎
The following theorem shows that $u^{2}(k)$ is exactly the optimal controller
for centralized information structure with two step delay, where the
information set of each vehicle is
$\mathcal{I}_{i}(k)=\left\\{x(0:k-2)\right\\}$.
###### Theorem 2
_Given that Assumption $1$ holds, an optimal solution to (16) is given by_
$\displaystyle u^{2}(k)=L(k)x^{2}(k)$ (18)
###### Proof:
Consider the system (17) together with the cost function $J^{2}_{u}$. Both
$x(0:k-2)$ and $u^{2}(k)$ are linear functions of $x(0:k-2)$ which is
independent of $w(k-2)$. Hence, finding the optimal control $u^{2}(k)$ is now
a centralized LQR problem. Applying proposition 2, we obtain (18). ∎
We now turn to the optimization problem (15). Recalling the expansions of
$x^{1}(k)$ and $u^{1}(k)$ in terms of $w(k-1)$ and $w(k-2)$, the expected
value of the $k^{th}$ term of $J^{1}_{u}$ can be expanded as follows:
$\displaystyle\textbf{E}\\{$
$\displaystyle\bigl{(}u^{1}(k)-L(k)x^{1}(k)\bigr{)}^{T}H(k)\bigl{(}u^{1}(k)-L(k)x^{1}(k)\bigr{)}\\}$
$\displaystyle=\textbf{Tr}\\{H(k)(F(k)-L(k))W(F(k)-L(k))^{T}\\}$
$\displaystyle\hskip
8.53581pt+\textbf{Tr}\\{H(k)\bigl{(}M(k)-L(k)(A+BF(k-1))\bigr{)}W$
$\displaystyle\hskip
68.28644pt\times\bigl{(}M(k)-L(k)(A+BF(k-1))\bigr{)}^{T}\\}$
where we used Proposition $4$a in the appendix and the fact that $w(k-1)$ and
$w(k-2)$ are independent. To minimize $J^{1}_{u}$ with respect to
$F(0),\ldots,F(k)$ and $M(1),\ldots,M(k)$, the difficulty is that $F$ and $M$
must satisfy specified sparsity constraints. We use vectorization of matrices
to simplify our optimization problem.
Let us define the matrix $D(k)$ as follows
$\displaystyle
D(k)\triangleq\begin{bmatrix}F(k-1)&M(k)\end{bmatrix}\;\;\in\mathbb{R}^{m\times
2p},\;k=1,\ldots,N-1$
and $D(N)\triangleq F(N-1)$. Then $\textrm{vec}\bigl{(}D(k)\bigr{)}$ is given
by
$\begin{bmatrix}\textrm{vec}\bigl{(}F(k-1)\bigr{)}\\\
\textrm{vec}\bigl{(}M(k)\bigr{)}\end{bmatrix}\;\;\;\in\mathbb{R}^{2mp},\;k=1,\ldots,N-1$
and $\textrm{vec}(D(N)\bigr{)}=\textrm{vec}\bigl{(}F(N-1)\bigr{)}$. Because of
the specified sparsity of $F$ and $M$, some components of
$\textrm{vec}\bigl{(}D(k)\bigr{)}$ must be zero. Let $S$ be the index set of
non-zero elements of $\textrm{vec}\bigl{(}D(k)\bigr{)}$, i.e.
$S\triangleq\left\\{i:\;\textrm{vec}_{i}\bigl{(}D(k)\bigr{)}\neq 0\right\\}$
Note that $\textrm{vec}\bigl{(}D(k)\bigr{)}$ and
$\textrm{vec}^{\star}\bigl{(}D(k)\bigr{)}$ are related by nonsquare matrix. We
define this matrix to be $E$, where dimensions implied by the context, so that
$\textrm{vec}\bigl{(}D(k)\bigr{)}=E\textrm{vec}^{\star}\bigl{(}D(k)\bigr{)}$.
The columns of $E$ are ${e_{j}}$ for ${j\in S}$ where $e_{j}$ denotes a column
vector having all zeros except a $1$ at the $j^{th}$ position. Since exactly
one entry in each column of $E$ is equal to $1$, $E^{T}XE$ is a sub-matrix of
$X$ containing exactly those rows and columns corresponding to the set $S$. We
illustrate the above definition via an example. Let
$D=\text{diag}(d_{11},d_{22},d_{33})\in\mathbb{R}^{3\times 3}$. For this
matrix, $S=\\{1,5,9\\}$, $E=\begin{bmatrix}e_{1}&e_{5}&e_{9}\end{bmatrix}$,
and $\textrm{vec}^{\star}(D)=[d_{11},d_{22},d_{33}]^{T}$.
In the following lemma, we show that a vectorization of matrices $F$ and $M$
makes the cost function $J^{1}_{u}$ a sum of quadratic functions without
constraints.
###### Lemma 4
_Define_
$\displaystyle Y_{11}(k)$
$\displaystyle=W\otimes(H(k-1)+B^{T}L^{T}(k)H(k)L(k)B)$ $\displaystyle
Y_{12}(k)$ $\displaystyle=-W\otimes B^{T}L^{T}(k)H(k)$ $\displaystyle
Y_{22}(k)$ $\displaystyle=W\otimes H(k)$
_and let_
$\displaystyle Y_{k}$ $\displaystyle=\begin{bmatrix}Y_{11}(k)&Y_{12}(k)\\\
Y_{12}^{T}(k)&Y_{22}(k)\end{bmatrix}$ $\displaystyle b_{k}$
$\displaystyle=\begin{bmatrix}Y_{22}(k-1)\\\
0\end{bmatrix}\textrm{vec}\left(L(k-1)\right)+\begin{bmatrix}Y_{12}(k)\\\
Y_{22}(k)\end{bmatrix}\textrm{vec}\left(L(k)A\right)$
_for $k=1,\ldots,N-1$, and_
$\displaystyle Y_{N}$ $\displaystyle=W\otimes H(N-1)$ $\displaystyle b_{N}$
$\displaystyle=\bigl{(}W\otimes H(N-1)\bigr{)}\textrm{vec}\left(L(N-1)\right)$
_Then optimization problem ( 15) is equivalent to_
$\displaystyle\min_{\textrm{vec}^{\star}(D(k))}\sum_{k=1}^{N}$
$\displaystyle\frac{1}{2}\textrm{vec}^{\star}(D(k))^{T}[Y_{k}]_{SS}\textrm{vec}^{\star}(D(k))$
$\displaystyle-\textrm{vec}^{\star}(D(k))^{T}[b_{k}]_{S}$
_Moreover, $Y_{k}$ is positive definite._
###### Proof:
See appendix. ∎
The advantage of this equivalent reformulation of the problem is that we have
$N$ quadratic functions without constraints and thus the optimal controller
gains can be computed by simply minimizing these functions separately.
###### Theorem 3
_Suppose $W$ is positive definite and Assumption $1$ holds. Then the optimal
gains of controllers are given by:_
$\displaystyle\textrm{vec}^{\star}(F(k-1))$
$\displaystyle=\begin{bmatrix}I&0\end{bmatrix}\textrm{vec}^{\star}(D(k))$
$\displaystyle\textrm{vec}^{\star}(M(k))$
$\displaystyle=\begin{bmatrix}0&I\end{bmatrix}\textrm{vec}^{\star}(D(k))$
_for $k=1,\ldots,N-1$ and
_$\textrm{vec}^{\star}(F(N-1))=\textrm{vec}^{\star}(D(N))$, _where_
$\textrm{vec}^{\star}(D(k))=[Y_{k}]_{SS}^{-1}[b_{k}]_{S}$.
### IV-D Steady State Controller Derivation
Assume that the solution to algebraic Riccati equation (11), $X(k)$, converges
to the stabilizing solution as $k$ approaches $\infty$:
$\displaystyle X=A^{T}XA+Q+A^{T}XB(B^{T}XB+R)^{-1}B^{T}XA$
Since $H(k)$ and $L(k)$ are specified by $X(k)$, they respectively converge to
matrices $H$ and $L$ as follows:
$\displaystyle H=B^{T}XB+R,~{}L=(B^{T}XB+R)^{-1}B^{T}XA$
Then $Y_{k}$ and $b_{k}$ will approach the values of $Y$ and $b$ given by
$\displaystyle Y$
$\displaystyle=\begin{bmatrix}W\otimes(H+B^{T}L^{T}HLB)&-W\otimes
B^{T}L^{T}H\\\ -W\otimes HLB&W\otimes H\end{bmatrix}$ $\displaystyle b$
$\displaystyle=\begin{bmatrix}W\otimes H\\\
0\end{bmatrix}\textrm{vec}(L)+\begin{bmatrix}-W\otimes B^{T}L^{T}H\\\ W\otimes
H\end{bmatrix}\textrm{vec}(LA)$
Thus, the optimal gains are calculated to be
$\displaystyle\textrm{vec}^{\star}(F)=\begin{bmatrix}I&0\end{bmatrix}[Y]_{SS}^{-1}[b]_{S}$
$\displaystyle\textrm{vec}^{\star}(M)=\begin{bmatrix}0&I\end{bmatrix}[Y]_{SS}^{-1}[b]_{S}$
### IV-E Estimation Structure
Having determined the optimal controller, we turn now to analyze this result.
Define $\zeta(k)=x(k)-w(k-1)$. Hence, we obtain the following state-space
system
$\zeta(k+1)=Ax(k)+Bu(k)$
with initial condition $\zeta(0)=0$. Note that the assumptions about the
information structure and sparsity structure of $A$ and $B$ guarantee that
each vehicle can update $\zeta(k)$ at time $k$. For example, consider Vehicle
$1$. Since Vehicle $1$ has access to $x_{2}(k-1)$ at time $k$, It can
construct $\zeta_{1}(k)=A_{11}x_{1}(k-1)+A_{12}x_{2}(k-1)+B_{1}u_{1}(k-1)$.
Letting $\xi(k)=\textbf{E}\\{x(k)|x(0:k-2)\\}$ the optimal control policy can
be written as
$\displaystyle u(k)=F(x(k)-\zeta(k))+M(x(k-1)-\zeta(k-1))+L\xi(k)$
In order to fully specify $u(k)$, the conditional estimates $\xi(k)$, as well
as the matrices $L$, $F$ and $G$ must be computed. We have
$\displaystyle\xi(k+1)$ $\displaystyle=\textbf{E}\\{x(k+1)\mid x(0:k-1)\\}$
$\displaystyle=A\textbf{E}\\{x(k)\mid x(0:k-1)\\}+B\textbf{E}\\{u(k)\mid
x(0:k-1)\\}$ $\displaystyle=A\zeta(k)+BM(x(k-1)-\zeta(k-1))+BL\xi(k)$
## V Numerical Results
We evaluate the performance of the system with the controller by giving an
example of a realistic scenario that HDV platoons often face on the road. In
practice, varying traffic conditions often mandate a deviation in the lead
vehicle’s velocity. Therefore, integral action for the lead vehicle is added
as a state to the system presented in (2), to model such disturbances.
Figure 2: Three HDV platoon, where a disturbance in velocity of the lead
vehicle is imposed. The top plot shows the velocity trajectories for the $M=3$
HDV platoon, the bottom plot shows the intermediate spacings, and the bottom
plot shows the control inputs. The trajectories are obtained through the
optimal distributed control and subindexed $i$, where $i=1,2,3$ denote the
platoon position index.
We consider a heterogeneous platoon, where the masses are set to
$[m_{1},m_{2},m_{3}]=[30000,40000,30000]$ kg. All the vehicles are assumed to
be traveling in the steady state velocity $v_{0}=19.44$ m/s ($70$ km/h) at
time gap $\tau=0.25$ s, which gives an intermediate distance of $d_{0}=4.86$.
The maximum engine and braking torque for a commercial HDV varies based upon
vehicle configuration but can be approximated to be 2500 Nm and 60000 Nm/Axle
respectively.
State disturbances as well as several lead vehicle deviation disturbances are
imposed on the system (Fig. 2). The lead vehicle deviation disturbances can be
explained by the following scenario. The platoon travels along a road where
the road speed is 70 km/h. Suddenly a slower vehicle enters the lane through a
shoulder path at a lower speed. The lead vehicle must therefore reduce its
speed to 60 km/h, in turn forcing the follower vehicles to reduce their speed
and adapt their relative distance accordingly. After a while, the slower
vehicle increases its speed to the road speed of 70 km/h and no longer
inhibits the platoon. Hence, the lead vehicle again resumes the road speed and
the follower vehicles adapt the speed and distance automatically as well.
Finally, the platoon arrives at a point where the road speed is changed to 80
km/h.
Fig. 2 shows the velocity trajectories of three HDV platoon in the top plot,
the corresponding intermediate spacings in the middle plot, and the required
control input to handle the disturbances in the bottom plot. The trajectories
nearly lie on top of each other, showing that the proposed controller performs
a tight control and the disturbances are handled well. There is no overshoot
in the velocity or intermediate spacing tracking. Furthermore, the control
input is well within the feasible physical range. The weight normalized
control input energy required to handle the imposed disturbances is reduced by
15 % for the second vehicle and 14 % for the third vehicle, with respect to
the first vehicle. Hence, the controller displays a fuel efficient behavior,
since the input energy is directly proportional to the fuel consumption. The
theoretical value, in this case, for the cost function with the proposed
optimal distributed control is only 0.01 % higher than a fully centralized
control with full state information at all times. On the other hand, the
proposed controller produces a 67 % lower theoretical cost compared to a
centralized control with two step time delays.
## VI SUMMARY AND CONCLUSIONS
We have presented an analytical controller, which is optimal under a delayed
information sharing pattern for chain structures. A discrete time HDV platoon
model has been derived that includes physical coupling with both neighboring
vehicles. The results show that the cost function with proposed controller is
very close to the fully centralized cost and better than the cost for the
centralized case with two time delays. Hence, the cost function can be
significantly reduced by considering additional available local information.
The controller maintains a tight control even though time delays are imposed.
For future work, we plan to extend to the presented results to $M$ HDVs and
arbitrary time step delays, which is relevant for HDV platooning.
## References
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### -A Preliminaries
###### Proposition 3 ([21])
_If $A$, $B$, $C$, $D$ and $X$ are suitably dimensioned matrices, then_
1. a)
$\textrm{vec}(AXB)=(B^{T}\otimes A)\textrm{vec}(X)$,
2. b)
$(A\otimes B)(C\otimes D)=(AC)\otimes(BD)$,
3. c)
_If $A$ and $B$ are positive definite, then so is $A\otimes B$,_
4. d)
$\textbf{Tr}\\{AXBX^{T}\\}=\textrm{vec}^{T}(X)(B^{T}\otimes
A)\textrm{vec}(X)$,
5. e)
$(A\otimes B)^{-1}=A^{-1}\otimes B^{-1}$._( $A$ and $B$ are nonsingular)_
###### Proposition 4 ([20])
_Let $x$ and $y$ be zero-mean random vectors with a jointly Gaussian
distribution. Assume $S$ be a symmetric matrix. Then the following facts
hold:_
1. a)
$\textbf{E}\\{x^{T}Sx\\}=\textbf{Tr}\bigl{\\{}S\textbf{{E}}\\{xx^{T}\\}\bigr{\\}}$.
2. b)
$\textbf{E}\\{x|y\\}$ and $x-\textbf{E}\\{x|y\\}$ are independent.
###### Proposition 5 ([21])
_Suppose that a symmetric matrix is partitioned as $\begin{bmatrix}A&B\\\
B^{T}&C\end{bmatrix}$, where $A$ and $C$ are square. This matrix is positive
definite if and only if $C$ and $\bigtriangleup=A-BC^{-1}B^{T}$ are positive
definite._
### -B Proof Lemma 3
First note that $x^{2}(k)=x(k)-x^{1}(k)$. Thus,
$\displaystyle x^{2}(k+1)=$ $\displaystyle
Ax(k)+Bu(k)-\bigl{(}A+BF(k)\bigr{)}w(k-1)$ $\displaystyle=$ $\displaystyle
Ax^{1}(k)+Ax^{2}(k)+Bu^{1}(k)+Bu^{2}(k)$
$\displaystyle-\bigl{(}A+BF(k)\bigr{)}w(k-1)$
The proof is completed by substituting
$x^{1}(k)=w(k-1)+\bigl{(}A+BF(k-1)\bigr{)}w(k-2)$ and
$u^{1}(k)=F(k)w(k-1)+M(k)w(k-2)$ into the second line.
### -C Proof Lemma 4
The equivalence of optimization problems (15) and (4) follows simply by
vectorization of matrices. First note that
$\textrm{vec}\bigl{(}F(k-1)\bigr{)}=\begin{bmatrix}I&0\end{bmatrix}\textrm{vec}\bigl{(}D(k)\bigr{)}$.
Thus
$\displaystyle\textrm{vec}(F^{-}-L^{-})$
$\displaystyle=[I\;0]\textrm{vec}\bigl{(}D(k)\bigr{)}-\textrm{vec}(L^{-})$
From Propositions $3$b and $3$d, we have
$\displaystyle\textbf{Tr}\left\\{H^{-}(F^{-}-L^{-})W(F^{-}-L^{-})^{T}\right\\}$
$\displaystyle=$
$\displaystyle\textrm{vec}^{T}\bigl{(}D(k)\bigr{)}\begin{bmatrix}W\otimes
H^{-}&0\\\ 0&0\end{bmatrix}\textrm{vec}\bigl{(}D(k)\bigr{)}$
$\displaystyle-2\textrm{vec}^{T}(L^{-})\begin{bmatrix}W\otimes
H^{-}&0\end{bmatrix}\textrm{vec}\bigl{(}D(k)\bigr{)}$
$\displaystyle+\textrm{vec}^{T}(L^{-})(W\otimes H^{-})\textrm{vec}(L^{-})$
Likewise,
$\textrm{vec}\bigl{(}M(k)\bigr{)}=\begin{bmatrix}0&I\end{bmatrix}\textrm{vec}\bigl{(}D(k)\bigr{)}$.
Then
$\displaystyle\textrm{vec}\bigl{(}M(k)-L(A+BF^{-})\bigr{)}=$
$\displaystyle\hskip
56.9055pt\bigl{(}[0\;\;I]-[LB\;\;0]\bigr{)}\textrm{vec}\bigl{(}D(k))\bigr{)}-\textrm{vec}(LA)$
Therefore,
$\displaystyle\textbf{Tr}\bigl{\\{}H\bigl{(}M-L(A+BF^{-})\bigr{)}W\bigl{(}M-L(A+BF^{-})\bigr{)}^{T}\bigr{\\}}$
$\displaystyle=$ $\displaystyle\textrm{vec}^{T}(D)\begin{bmatrix}W\otimes
B^{T}L^{T}HLB&-W\otimes B^{T}L^{T}H\\\ -W\otimes HLB&W\otimes
H\end{bmatrix}\textrm{vec}(D)$
$\displaystyle-2\textrm{vec}^{T}(LA)\begin{bmatrix}-W\otimes HLB&W\otimes
H\end{bmatrix}\textrm{vec}\bigl{(}D(k)\bigr{)}$
$\displaystyle+\textrm{vec}^{T}(LA)(W\otimes H)\textrm{vec}(LA)$
After Substituting these values back into $J_{u}^{1}$, using
$\textrm{vec}(D)=E\textrm{vec}^{\star}(D)$, and eliminating constant terms, we
arrive at (4).
The only part that remains to be proved is that $Y_{k}$ is positive definite.
Since $W$ and $H(k)$ are positive definite, $Y_{22}(k)$ is positive definite
according to Proposition $3$c. Proposition $3$e then implies that
$Y_{22}^{-1}(k)=W^{-1}\otimes H^{-1}(k)$. From Proposition $3$b, we have
$\displaystyle Y_{12}(k)Y_{22}^{-1}(k)Y^{T}_{12}(k)=W\otimes
B^{T}L^{T}(k)H(k)L(k)B$
Consequently,
$\displaystyle\bigtriangleup(k)$
$\displaystyle=Y_{11}(k)-Y_{12}(k)Y_{22}^{-1}(k)Y^{T}_{12}(k)$
$\displaystyle=W\otimes H(k-1)$
Since $\bigtriangleup(k)$ and $Y_{22}(k)$ are positive definite, from
Proposition $5$, $Y_{k}$ is positive definite. Finally note that $E^{T}$ has
full row rank, so $[Y]_{SS}=E^{T}Y_{k}E$ is positive definite.
|
arxiv-papers
| 2013-09-17T10:36:45 |
2024-09-04T02:49:51.001673
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Hamid Reza Feyzmahdavian, Assad Alam, Ather Gattami",
"submitter": "Hamid Reza Feyzmahdavian",
"url": "https://arxiv.org/abs/1309.4251"
}
|
1309.4290
|
Simple free-surface detection in two and three-dimensional SPH solver
Agra Barecasco1,2, Hanifa Terissa 1,2, Christian Fredy Naa2
1Graduate School of Natural Science and Technology, Kanazawa University
Kakuma, Kanazawa 920-1192 Japan
2Faculty of Mathematics and Natural Sciences, Institut Teknologi Bandung
Jl. Ganesha 10, Bandung 40132 Indonesia
E-mail: [email protected], [email protected],
[email protected]
Abstract. A simple free-surface particle detection method for two and three-
dimensional SPH simulation has been implemented. The method uses sphere
representation for the SPH particle. The fluid domain is covered by
overlapping spheres. A sphere whose surface is not fully covered considered as
boundary. To test particle’s boundary status, we used a sum of normalized
relative position vectors from neighbouring particles to the test particle. By
checking the existence of uncovered sphere surface by this vector sum,
boundary status of the test particle can be determined. This boundary
detection method can be easily embedded in the SPH solver algorithm.
Keywords: Smoothed-particle Hydrodynamics, boundary particle detection
## 1 Introduction
In recent years there has been much development of smoothed-particle
hydrodynamics (SPH)[1]. Implementation of boundary conditions is not so clear
with SPH method compared with mesh based methods. In SPH, the problem that we
must first locate the points that are regarded as boundary.
For boundary particle detection, Randles and Libersky[1] previously have
suggested using the sums of the gradients of the SPH kernels. Ideally, these
kernel gradients sum to zero for interior particles. Any particle for which
the sum of the gradient kernels is not near zero is presumably an exterior
particle. This method gives correct results when the particles are uniformly
spaced. Aamer and Dilts[2] developed an algorithm using overlapping spheres as
fluid particles representation. The method can geometrically detect free-
surface particles in a robust way. However, its extension to three dimensional
simulations requires a large amount of calculation. Marrone[3] gave the idea
of using a scan cone around the expected normal vector of the fluid surface to
make a further check if there is any particle covering the test particle.
In this paper, we propose a simple boundary detection method using both the
idea of Dilts[2] and Marrone[3]. The aim is to detect uncovered sphere segment
on the boundary particle sphere. In our method, we develop a very simple
algorithm to estimate the existence of this uncovered sphere segment.
## 2 Numerical Schemes: Smoothed Particle Hydrodynamics
Smoothed-particle hydrodynamics is a method to simulate fluids[1]. It is a
lagrangian meshless method based on convolution of smoothing function (kernel)
$W$ over fluid’s field functions to approximate their values at a point. This
field function is discretized by a set of point particles. The kernel
approximation is given by the form
$f(\mathbf{x})=\int\limits_{\Omega}f(\mathbf{x}^{\prime})W(\mathbf{x}-\mathbf{x}^{\prime})d\mathbf{x}^{\prime},$
where $\Omega$ represents the convolution domain. In our simulation we use the
cubic spline kernel:
$W(\mathbf{x},h)=\beta\begin{cases}\frac{2}{3}-q^{2}+\frac{1}{2}q^{3},&0\leq
q<1\\\ \frac{1}{6}\left(2-q^{3}\right),&1\leq q<2\\\ 0,&q\geq 2,\end{cases}$
where $2h$ is the radius of the kernel function’s support and
$q=\frac{|\mathbf{x}|}{h}$.
To simulate an inviscid fluid, the convolution of this kernel is applied to
the following governing Euler equations for inviscid fluid motion:
$\displaystyle\frac{\mbox{D}\rho}{\mbox{D}t}$
$\displaystyle=-\rho\nabla\cdot\mathbf{v}$ (1)
$\displaystyle\frac{\mbox{D}\mathbf{v}}{\mbox{D}t}$
$\displaystyle=-\frac{1}{\rho}\nabla p+\mathbf{g}$ (2)
where $\mbox{D}/\mbox{D}t$ represent material derivative following an
infinitesemal fluid element. $\rho,\mathbf{v},p$ represent density, velocity
and pressure, respectively. The external acceleration $\mathbf{g}$ is given to
simulate gravity.
For SPH simulation, the fluid domain is discretized by fluid particles. The
discretized SPH approximation form of Euler equations above are given by
$\displaystyle\frac{\mbox{D}\rho_{i}}{\mbox{D}t}$ $\displaystyle=$
$\displaystyle\sum\limits_{j}m_{j}\left(\mathbf{v}_{i}-\mathbf{v}_{j}\right)\cdot\nabla
W(\mathbf{x}_{i}-\mathbf{x}_{j},h)$ (3)
$\displaystyle\frac{\mbox{D}\mathbf{v}_{i}}{\mbox{D}t}$ $\displaystyle=$
$\displaystyle\sum\limits_{j}m_{j}\left(\frac{p_{i}+p_{j}}{\rho_{i}\rho_{j}}\right)\nabla
W(\mathbf{x}_{i}-\mathbf{x}_{j},h)+\mathbf{g},$ (4)
where $i$ and $j$ denote particle indices and $m_{j}$ represent $j$-th
particle mass.
To determine the pressure, we used the state equation
$p_{i}=c^{2}(\rho_{i}-\rho_{0}),$ (5)
where $c$ is the speed of sound constant and $\rho_{0}$ is the reference
density. Here $\mathbf{x}$ represents particle position and $j$ sums to all
particles within the kernel radius. In the implementation, the sums must check
for all particles in the fluid domain if they are within the kernel. To make a
faster calculation, we used a linked-list grid method[5] to reduce the
calculation amount per particle so $j$ only runs through the nearest
neighbouring particles.
### SPH solver algorithm
Implementation of eq. (3) and eq. (4) involves many steps. Here we give the
general steps of SPH solver that are related to our boundary detection method.
For our cases, we used Leap-Frog integrator to advance in time. SPH simulation
can be run as follows:
1. 1.
Initialization. Set the initial condition of attributes of all SPH particles
2. 2.
Do the followings for all fluid particles until the run-time limit has been
reached:
1. (a)
Calculate the current pressure of all particles by using state equation eq.
(5)
2. (b)
For each particle, insert it into the linked-list grid[5] and register its
nearest neighbour particles
3. (c)
For each particle, calculate its momentum change rate (eq. (4)) by using the
new calculated pressure, then calculate its new velocity and position
4. (d)
For each particle, calculate its density change rate (eq. (3)) by using new
calculated velocity, then calculate its new density
We see that SPH simulation involves the calculation of summation from all
neighbouring particles to determine each particle’s density, pressure and
momentum. This is an advantage for our detection method which is explained in
the next section.
## 3 Boundary Particle Detection Method
### 3.1 Boundary particle definition
We are given a set of spheres in $\mathbb{R}^{3}$,
$S=\\{s_{1},s_{2},\dots,s_{n}\\}$. Each sphere $s_{i}$ represents a SPH
particle centered at the particle’s position $\mathbf{x}_{i}$. Let $r_{i}$ be
the sphere’s radius.
A point $\mathbf{x}\in\mathbb{R}^{3}$ is said to be _covered_ by $s_{i}$ if
$|\mathbf{x}-\mathbf{x}_{i}|<r_{i}$. SPH _fluid domain_
$F_{S}\in\mathbb{R}^{3}$ is a domain where every $\mathbf{x}\in F_{S}$ is
covered by one or more SPH spheres of $S$. $F_{S}$ is divided into a number of
subspaces by the sphere set $S$. The spheres in $F_{S}$ overlap each other
representing the compact fluid domain.
Consider two overlapping spheres $s_{i}$ and $s_{j}$. A _spherical cap_
$Cap(i,j)$ of $s_{j}$ on $s_{i}$ is a spherical segment of $s_{i}$’s surface
where every point on it is covered by $s_{j}$. Every $Cap(i,j)$ has to be
bounded by an intersection circle $Circ(i,j)$ on the surface of $s_{i}$ (see
figure 1). The $Circ(i,j)$ is said to be _covered_ if all points in it are
covered.
Figure 1: Illustration of terminologies
Now we consider some special cases. The radius of all spheres in $S$ is the
same, so the situation that a sphere is located inside a sphere is impossible.
To avoid two spheres or more to coincide, we improved the SPH algorithm with
artificial viscousity[1] so that there is no interparticle penetration.
Definition 3.1. Let sphere $s_{i}$ be overlapped by $n$ spheres. The sphere
$s_{i}$ is said to be _covered_ if every point on sphere $s_{i}$ is covered by
at least one of the $n$ intersecting spheres. The particle of the
corresponding sphere is then said to be an _interior_ SPH particle.
Figure 2: Illustration of a boundary particle
Definition 3.2. Let sphere $s_{i}$ be overlapped by $n$ spheres. The
corresponding SPH particle is said to be a _boundary_ particle if there exists
a point on the corresponding sphere which is not covered by any of the
overlapping spheres (see figure 2).
The smoothness of the fluid surface depends on the value of $r_{i}$. If we use
the spheres of smaller size, the boundary particles have greater resolution,
which gives a sharper surface profile (see figure 3). However, a small sized
sphere representation bring the risk of geometrical cavity to emerge inside
the fluid. This could lead to an error of boundary detection.
Figure 3: Drawn in 2D, the two figures show how a smaller representative
sphere radius give sharper surface profile. Uncovered sphere segments are
shown in red.
### 3.2 Detection method
Based on the definitions above, now we explain our method to detect boundary
particles.
Let $s_{i}$ be a SPH particle represented as a sphere with particle’s position
$\mathbf{x}_{i}$ as the center. Within its kernel support, $m$ neighbour
particles are present to contribute for $s_{i}$’s SPH attributes. Let these
neighbouring particle spheres overlap $s_{i}$. The _cover vector_
$\mathbf{b}_{i}$ is defined as (see figure 4):
$\mathbf{b}_{i}=\sum\limits_{j=0}^{n}\frac{\mathbf{x}_{i}-\mathbf{x}_{j}}{|\mathbf{x}_{i}-\mathbf{x}_{j}|}$
(6)
Let $l_{i}$ be a ray that starts from $\mathbf{x}_{i}$ and parallel to
$\mathbf{b}_{i}$. If $l_{i}$ does not intersect any spherical cap, then the
corresponding sphere is not covered. Boundary particle detection can be done
by checking the existence of $l_{i}-Cap$ intersection.
Figure 4: Illustration of cover vector and the scan cone, drawn in 2D. The
cone checks if there is any other covering particle within kernel radius.
Now we consider the case for particle located at or near the boundary. Because
of the non-uniform property of the particles (see figure 5), the direction of
$\mathbf{b}_{i}$ only serves as a rough estimation of the normal of the
surface. We used a _scan cone_ (see figure 4) which checks if there is any
covering particle near $l_{i}$. The cone angle $\theta_{i}$ serves as
threshold angle to any particle that lies near $l_{i}$, which detects the
existence of a spherical cap intersection with $l_{i}$.
Figure 5: The cover vector on each particle differs according to the
distribution of neighbouring particles.
The value of $\theta_{i}$ depends on the radius of representative sphere used
and plays important role in checking the existence of $l_{i}-Cap$
intersection. In our case, we directly chose the value $\theta_{i}=\pi/3$. The
scan cone calculation check for particle $i$ is given by
$\mbox{IF
}\arccos\left(\frac{\mathbf{x}_{j}-\mathbf{x}_{i}}{|\mathbf{x}_{j}-\mathbf{x}_{i}|}\cdot\frac{\mathbf{b}_{i}}{|\mathbf{b}_{i}|}\right)\leq\frac{\theta_{i}}{2}\mbox{
THEN assign boundary status to the particle $i$}$ (7)
In case when the fluid particles join to form a thin jet, a plane or become
solitary particle cluster the number of neighbouring particles for each
becomes very small. We considered them as boundary particles by giving a
threshold number
$n_{th}=\begin{cases}4&\mbox{ for 2 dimensional case}\\\ 15&\mbox{ for 3
dimensional case}\end{cases}$ (8)
to the number of neighbour particles. A particle having the number of
neighbour particles greater than $n_{th}$ must be tested for its boundary
status, otherwise it is a boundary particle.
### 3.3 Boundary detection algorithm
Our algorithm is embedded within the SPH solver algorithm that has been
mentioned in section 2 above. The summation to calculate the cover vector can
be done simultaneously with any SPH solver step that involves summation
through all neighbouring particles. The simple surface detection algorithm
written below uses some steps of SPH solver algorithm simultaneously:
1. 1.
At _initialization_ , assign all particles in the fluid domain as boundary
particles
2. 2.
During the step of registering particles into linked-list grid, for each
particle $i$ calculate the cover vector (eq. (6)) by summing
$\frac{\mathbf{x}_{i}-\mathbf{x}_{j}}{|\mathbf{x}_{i}-\mathbf{x}_{j}|}$
for every detected neighbour particle $j$
3. 3.
During the step of calculating momentum change rate for all particles (eq.
(4)), for each particle $i$, if the number of neighbours of $i$ is smaller
than or equal to the threshold $n_{th}$ (eq. (8)) then consider the particle
$i$ as boundary
4. 4.
During the step of calculating density change rate for all particles (eq.
(3)), check the existence of $l_{i}-Cap$ intersection by applying the scan
cone calculation (eq. (7)). If the intersection exists, assign the current
$i$-th particle as interior.
## 4 Test Cases
We have applied the detection algorithm to several cases. To visualize
particles we used spheres of radius $0.3h_{sm}$, where $h_{sm}$ is the kernel
support radius. Red and blue color was used to distinguish between boundary
and interior particles.
Figure 6: Two dimensional dam break
### 4.1 Two-dimensional dam break
We implemented the algorithms to two-dimensional case so the result can be
seen clearly. The simulation is a dam break in rectangular tank. The
simulation consists of 2997 SPH particles with the initial shape as a
rectangular block of fluid at the right side of the tank. Figure 6 shows the
result.
### 4.2 Two-dimensional standing wave
The simulation is a standing wave in a rectangular tank with periodic boundary
condition at the sides. The simulation consists of 3720 SPH particles with the
initial condition using a sine function. Figure 7 shows the result.
Figure 7: Two dimensional standing wave
Figure 8: Three dimensional droplet motion.
Figure 9: Three-dimensional standing wave. The upper left figure shows the
system before boundary detection.
### 4.3 Three dimensional droplet motion
This simulation applied surface tension model used by Zhang[4] to the boundary
particles. The droplet consists of 3375 fluid particles. The surface tension
depends on the curvature of the surface shape that calculated by applying
Moving Least Square fit to the boundary particles as interpolation points.
Figure 8 shows how surface tension deforms the initial cube shape into a
sphere.
### 4.4 Three-dimensional standing wave
This case is a 3-dimensional standing wave with periodic boundary condition at
the sides of a cube tank. We used 44191 particles. The initialization of the
dome-shaped surface is given by a gaussian function. Figure 9 shows the
result.
## 5 Correction Method: Scan Circle
The implementation of boundary detection method worked well in determining the
deep interior particles, but a drawback occured when determining the boundary
particles. Because of the unpredictable behavior of the neighbouring particle
distribution, the cover vector of a boundary particle sometimes aims too close
to another neighbour particle. The scan cone mistakenly regarded this boundary
particle as interior (see figure 10). This impairs the accuracy of this
method.
Figure 10: The scan cone of boundary particle $i$ regards another boundary
particle as a covering particle. Particle $i$ is mistakenly considered as
interior. Figure 11: Scan circle
To treat this problem, a further research for implementing a _scan circle_ is
on going. Consider an intersection given by the scan cone and the sphere of a
SPH particle $i$. The intersection is a circle whose center located in ray
$l_{i}$ (see figure 11). The existence of uncovered sphere segment can be
tested by checking whether this circle is covered or not.
## 6 Summary
The simple boundary detection in SPH based on the implementation of cover
vector and the scan cone has been successfully implemented to several cases.
The algorithm can be embedded in the SPH solver algorithm easily. However, a
drawback exists in the boundary detection accuracy because of random behavior
of particle distribution. A research on correction method using a scan circle
coverage is on going.
## References
* [1] P.W. Randles, L.D. Libersky (1996), Smoothed particle hydrodynamics: some recent improvements and applications, Comput. Meth. Appl. Mech. Eng. , 139, 375-408
* [2] H. Aamer, G.A. Dilts (2007), Three-dimensional boundary detection for particle methods, J. Comp. Phys. , 226, 1710-1730
* [3] S, Marrone et.al. (2010), Fast free-surface detection and level-set function definition on SPH solvers, J. Comp. Phys. , 229, 3652-3663
* [4] M. Zhang, (2010), Simulation of surface tension in 2D and 3D with smoothed particle hydrodynamics method, J. Comp. Phys. , 229, 7238-7259
* [5] G.R. Liu, M.B. Liu (1997). Smoothed Particle Hydrodynamics: A Meshfree Particle Method, World Scientific Publishing
|
arxiv-papers
| 2013-09-17T12:57:23 |
2024-09-04T02:49:51.011792
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Agra Barecasco, Hanifa Terissa, Christian Fredy Naa",
"submitter": "Christian Naa",
"url": "https://arxiv.org/abs/1309.4290"
}
|
1309.4294
|
# Projectile fragment emission in fragmentation of 56Fe on C, Al, and CH2
targets at 471 A MeV ††thanks: Submitted to Chin. Phys. C
Yan-Jing Lia, Dong-Hai Zhangb, Shiwei Yana,c , Li-Chun Wanga,
Jin-Xia Chenga, Jun-Sheng Lib, S. Kodairad, and N. Yasudad
aCollege of Nuclear Science and Technology, Beijing Normal University, Beijing
100875, China
bInstitute of Modern Physics, Shanxi Normal University, Linfen 041004, China
cBeijing Radiation Center, Beijing 100875, China
dFundamental Technology Center, National Institute of Radiological Sciences
4-9-1 Anagawa, Inage-ku, Chiba 263-8555, Japan Corresponding author. Tel:
+863572051347; fax: +863572051347. E-mail address:[email protected]
Abstract
The emission angle and the transverse momentum distributions of projectile
fragments produced in fragmentation of 56Fe on CH2, C, and Al targets at 471 A
MeV are measured. It is found that for the same target the average value and
width of angular distribution decrease with increase of the projectile
fragment charge, and for the same projectile fragment the average value of the
distribution increases and the width of the distribution decreases with
increasing the target charge number. The transverse momentum distribution of
projectile fragment can be explained by a single Gaussian distribution and the
averaged transverse momentum per nucleon decreases with the increase of the
charge of projectile fragment. The cumulated squared transverse momentum
distribution of projectile fragment can be well explained by a single Rayleigh
distribution. The temperature parameter of emission source of projectile
fragment, calculated from the cumulated squared transverse momentum
distribution, decreases with the increase of the size of projectile fragment.
PACS 25.70.-z, 25.70.Mn, 29.40.Wk
## 1 Introduction
The knowledge of heavy ion fragmentation at intermediate and high energy is
very important in nuclear physics, astrophysics, and medical physics.
Considering the biological effects of space radiation, when astronaut have
their mission outside the earth magnetic field, they are suffered from
Galactic Cosmic Radiation(GCR) and Solar particle events, e.g., showers of
energetic charged particles from the surface of the Sun. These energetic
charged particles will be the dominant sources of the radiation dose and
affect the health of humans on long-duration spaceflight both inside and
outside the station. According to the GCR model developed by Badhwar and
O’Neill[1], in unshielded free space in the inner heliosphere, iron ions
deliver about $8\%$ of the total dose from the GCR and $27\%$ of the dose
equivalent at times near solar maximum, even though they contribute less than
$1\%$ of the total GCR flux. Because iron ions are the most densely ionizing
particles which are presented in significant numbers in the GCR, there has
been considerable interest in understanding their transport through matter and
their biological effects.
The understanding of the radiobiology of heavy charged particles (HZE) is a
subject of great interest due to the complicated dependence of their relative
biological effectiveness on the type of ion and its energy, and its
interaction with various targets. It has become clear that heavy ions have the
largest radiological effects. These effects also appear in regions close to
the beam entrance, i.e., in the depth-dose plateau region, where normal tissue
is usually situated. In addition, due to the longer ranges of the fragments
produced by the fragmentation of the incident ions, the tails of the dose
distribution beyond the Bragg peak may be too high for minimizing doses to
normal tissue beyond the primary ion range. Finally, recent experimental
results[2] have revealed that the fragments are emitted at larger angles than
the scattering angles of the beam, which further increases the spread of the
beam. Exact information about the fragment emission angular distributions will
be especially important in radiotherapy. So far only a few measurements have
been performed to analyze fragment emission angles from HZE reactions below
500 A MeV[2-4].
Fragmentation is a term commonly used to specify a nuclear disassembly by
force. Hot fragmentation is meant to indicate the most violent of these
process, following excitation beyond the limits of nuclear binding, but still
ending with bound nuclear fragments of different sizes in the final
channels[5]. The formation mechanism of these fragments, whether they the
remnants of an incomplete destruction or the products of a condensation out of
the disordered matter, has continued to be the topic of very active research.
In order to describe the physical process of heavy ion transport, several one-
dimensional Monte Carlo codes, such as HZETRN[6], HIBRAC[7], FLUKA[8],
NUCFRAG2[9], and three dimensional model[10] are appeared. The improved
quantum molecular dynamics model (ImQMD) is a dynamical model which is
developed to follow the reaction process on a microscopic level[11, 12].
The properties of 56Fe on various targets at various energies have been
studied by many groups[13-25], most of the studies are devoted to the total
charge-changing cross sections and the partial cross sections of fragment
productions, a little attention is paid to the fragment emission angular
distribution study.
In this paper, we present the results of the emission angular distribution,
transverse momentum distribution and the temperature of emission source of
fragment produced in fragmentation of 471 A MeV 56Fe on C, Al, and CH2
targets. CH2 target is used to obtain the cross section on a hydrogen target.
The fragmentation cross sections is published in our previous paper[26]. The
arrangements of this paper are as follows: we introduce our experimental
detail in sec. II. In sec. III, experiment result and discussion are given. At
last, we give the conclusion in sec. IV.
## 2 Experimental details
### 2.1 Experimental setup
Stacks of C, Al and CH2 targets sandwiched with CR-39 detectors were exposed
normally to 471 A MeV 56Fe beams at the Heavy Ion Medical Accelerator in Chiba
(HIMAC) at the Japanese National Institute of Radiological Sciences (NIRS).
Figure 1 shows the configuration of sandwiched target. A CR-39 sheet, about
0.77 mm in thickness, is placed before and after the targets. The thickness of
carbon, aluminum and polyethylene targets is 5, 3, and 10 mm, respectively.
Figure 1: Sketch of the target-detector configuration.
After exposure, the CR-39 detectors are etched in 7N NaOH aqueous solution at
temperature of 70∘C for 15 hours. Then, the beam ions and their fragments
manifest in the CR-39 as etch-pit cones on both sides of CR-39 sheets. The
images of ion tracks are scanned and analyzed automatically by HSP-1000
microscope system and the PitFit track measurement software, then checked
manually. The PitFit software allows us to extract some geometric information,
such as the position coordinates, major and minor axes and area of etched
track spot on CR-39 surfaces. Image data ($45\times 45mm^{2}$) are acquired
for both front and back surfaces of each CR-39 detector. About 2$\times$104 Fe
ions are traced from the first CR-39 detector surface in the stack. 56Fe
trajectories and the ones of secondary fragments are reconstructed in the
whole stack.
### 2.2 Experimental method
The spots on the front surface (with respect to the beam direction) are
directly scanned firstly, then the CR-39 sheet is turned under the middle line
of the sheet and the spots on the back surface are scanned. The trajectories
of ion tracks through CR-39 sheets are reconstructed in two steps using the
track tracing method[27]: (1) the track position in CR-39 surfaces is
corrected by parallel and rotational coordinate translation (except for the
track position on upper surface of the first CR-39 sheet), and (2) the
difference between the track position of corresponding tracks on both side of
the CR-39 sheets and on the surfaces neighboring targets is minimized by a
track matching routine. The coordinate of track before the target (or front
surface of CR-39 sheet) is $(x,y)$ and of matching track after the target (or
back surface of CR-39 sheet) is $(x^{\prime},y^{\prime})$. Following the
translation relation, the coordinate of matching track can be calculated as:
$\displaystyle x^{\prime}_{th}=$ $\displaystyle ax+by+c,$ (1) $\displaystyle
y^{\prime}_{th}=$ $\displaystyle a^{\prime}x+b^{\prime}y+c^{\prime},$ (2)
parameters $a,b,c,a^{\prime},b^{\prime}$, and $c^{\prime}$ are determined
using the least square method. Then, the coordinate $x^{\prime}_{th}$,
$y^{\prime}_{th}$ of matching track is calculated. However, because of the
Coulomb scatter etc. contributions, $x^{\prime}_{th}$, $y^{\prime}_{th}$ are
certainly different from $x^{\prime},y^{\prime}$, the difference
$dx=x^{\prime}_{th}-x^{\prime}$, $dy=y^{\prime}_{th}-y^{\prime}$ is calculated
which can help us to determine the matching track.
Figure 2: The distribution of dx and dy. (a) and (b) the difference between
the front and back surface on a CR-39 sheet, (c) and (d) the difference before
and after carbon target. Figure 3: The area distribution of etch-pit spots of
all of 56Fe ions and ones of secondary fragments on CR-39 surface after CH2
target.
Figure 2 (a) and (b) show the difference $dx$ and $dy$ in the front surface
and back surface on a CR-39 sheet, (c) and (d) show the difference $dx$ and
$dy$ before and after the target. If the difference are calculated for all
combinations of positions for extracted tracks, only the matching combination
ought to make a peak which appears in the figures, and the difference $dx$ and
$dy$ of other combinations should be randomly distributed. The deviations
$\sigma(x^{\prime})$ and $\sigma(y^{\prime})$ give the position accuracies of
tracks which are estimated to be $2-4\mu$m for the case of (a) and (b), and
$8-20\mu$m for case of (c) and (d).
The matching iron ion track is searched within $4\times\sigma(x^{\prime})$ and
$4\times\sigma(y^{\prime})$ region of $x^{\prime}_{th}$ and $y^{\prime}_{th}$.
The matching projectile fragments are searched within the limited
fragmentation angle, which is about $10^{\circ}$ in present experiment. The
number of projectile fragments leaving the target is determined from the
distribution of the etched area. Figure 3 shows the track base area
distribution of 56Fe ions and their fragments in CR-39 sheet. Peaks for 56Fe
and its fragments with charge down to $Z=6$ appear clearly. Because of the
limitation of CR-39 detector, the tracks of fragments with charge
$1\leq{Z}\leq 5$ are not fully etched as a measurable spots. The emission
angle $\theta$ of each fragment is calculated by taking readings of the
coordinates of the beam track and the fragment track.
## 3 Results and discussion
Emission angular distribution and transverse momentum distribution of
projectile fragments provide information on the nuclear structure and the
mechanism of nuclear interaction through which fragments are produced. These
distributions are also very important in designing experiments with
radioactive nuclear beams.
### 3.1 Angular distribution
Emission angle of each fragment and scattering angle of iron ion is calculated
from the coordinates of track positions on the surface of CR-39 sheet after
the target. The angular uncertainty is determined using the quadruplet fitting
method[4]
$\displaystyle\sigma(\theta)=\frac{\sqrt{\sigma^{2}_{z}\sin^{2}\theta+2\sigma^{2}_{p}\cos^{2}\theta}}{2h},$
(3)
where $\sigma_{p}$ is the positional uncertainty in $x-y$ plane of the stack
coordinate system which is about $3\mu$m for C-target, $\sigma_{z}$ is the
positional uncertainty in the z-axis which comes from stack composition and
detector thickness measurement and is estimated at $\approx 8\mu$m, $\theta$
represents the emission angle of the fitted line. With a detector thickness of
$h\approx 780\mu$m we thus obtain angular uncertainty $\sigma(\theta)\approx
0.16^{\circ}$ for value of $\theta$ up to $8^{\circ}$.
Figure 4 shows the angular distribution of primary iron ions and their
fragments for different targets. The emission angle of primary iron ion mainly
comes from the Coulomb scattering, most of which is less than $0.6^{\circ}$.
Most of projectile fragments have the emission angle less than $1.5^{\circ}$,
but some of fragments have emission angle up to $8^{\circ}$. The position of
the maximum of the fragment angular distribution increases slightly with
increase of the target mass, which can be explained that with the increase of
target mass the interaction between projectile and target is increased and
influence of fragment suffered from target is increased.
Figure 4: The angular distribution of primary iron ions and their fragments
for different targets, for comparison the counts of fragments is enlarged two
times.
Figure 5: The emission angle distributions of fragments with charge Z=6, 18,
and 24, respectively.
Figure 6: The emission angle distributions of fragments with charge Z=24 and
25, respectively.
Fig. 5 shows the emission angle distributions of projectile fragments with
$Z=6,18$, and 24 produced from the fragmentation of 56Fe on C and CH2 targets,
and Fig. 6 shows the emission angle distributions of projectile fragments with
$Z=24$ and 25 produced from the fragmentation of 56Fe on C, Al and CH2 targets
respectively. For fragments from the fragmentation of 56Fe on C and Al
targets, each angular distribution is fitted by a single Gaussian
distribution, and for fragments from the fragmentation of 56Fe on CH2 target,
each angular distribution is fitted by two Gaussian distribution, the fitting
parameters including $\chi^{2}$/DOF are presented in Table 1, where DOF means
the degree of freedom of simulation. From Figs. 5 and 6 and results in Table 1
it shows that for the same target the average value and the width of the
distribution decrease with increasing the charge number of projectile
fragment, and for the same projectile fragment the average value of the
distribution increases and the width of the distribution decreases with
increasing the target charge number.
Table 1. Values of fitting parameters of angular distribution using Gaussian
distribution. Z Target percentage(1) mean value(1) error(1) percentage(2)
mean value(2) error(2) $\chi^{2}$/DOF 25 CH2 $16.38\pm 5.52$ $0.25\pm 0.02$
$0.13\pm 0.03$ $7.33\pm 2.92$ $0.49\pm 0.17$ $0.28\pm 0.08$ 1.12 C $20.89\pm
1.62$ $0.23\pm 0.02$ $0.21\pm 0.02$ 1.81 Al $33.28\pm 13.27$ $0.29\pm 0.02$
$0.10\pm 0.04$ 1.52 24 CH2 $17.73\pm 1.50$ $0.39\pm 0.02$ $0.21\pm 0.02$
$1.44\pm 0.54$ $1.13\pm 0.11$ $0.18\pm 0.16$ 1.15 C $14.66\pm 1.84$ $0.36\pm
0.03$ $0.26\pm 0.03$ 2.15 Al $21.68\pm 5.69$ $0.36\pm 0.05$ $0.22\pm 0.05$
0.31 18 CH2 $22.90\pm 3.61$ $0.53\pm 0.06$ $0.25\pm 0.04$ $10.93\pm 4.26$
$1.16\pm 0.08$ $0.18\pm 0.08$ 2.17 C $16.00\pm 3.58$ $0.62\pm 0.12$ $0.54\pm
0.14$ 0.40 6 CH2 $13.71\pm 0.95$ $1.86\pm 0.34$ $1.73\pm 0.31$ $0.30\pm 1.06$
$2.97\pm 1.07$ $0.04\pm 1.44$ 0.12 C $10.74\pm 3.66$ $2.04\pm 0.34$ $1.78\pm
0.81$ 1.01
### 3.2 Transverse momentum distribution
The transverse momentum per nucleon ($p_{t}$) of a projectile fragment was
calculated on the basis of its emission angle $\theta$,
$p_{t}=p\sin\theta,$ (4)
where p is the momentum per nucleon of beam which can be calculated from beam
energy per nucleon (E), $p=(E^{2}+2m_{0}E)^{1/2}$. $m_{0}$ is the nucleon rest
mass and $\theta$ the emission angle of the projectile fragment with respect
to the beam direction.
Figure 7 shows the transverse momentum distribution of projectile fragment
with charge $Z=6,18$, and 24 produced from the fragmentation of 56Fe on C and
CH2 targets. The distributions can be well fitted by a single Gaussian
distribution, which is the same as the distributions of fragments produced in
reactions of light projectiles[28, 29] and heavy projectiles[30, 31, 32, 33,
34] at relativistic energies. These Gaussians are in good agreement with
predictions of the statistical model of Goldhaber[35]. This model assumes that
the Fermi momenta of the nucleons in a fragment are statistically distributed
as those in the original projectile nucleus. The averaged transverse momentum
per nucleon for fragment with charge $Z=6,18$, and 24 are $35.89\pm 35.25$,
$11.34\pm 8.74$, and $6.91\pm 4.74$ A MeV/c, respectively, for C target, and
$34.38\pm 30.79$, $11.98\pm 7.73$, and $7.08\pm 4.25$ A MeV/c, respectively,
for CH2 target. For the same target the averaged transverse momentum per
nucleon, and the width of the distribution increase with the decrease of the
charge of fragment. This tendency is also observed in Ref.[34].
Projectile fragments come from the directly produced fragments (primary
fragments) and the sequential decay fragments from excited primary fragment.
However, since the primary fragments are excited, they are deexcited by light
particle evaporation. This secondary decay decreases the observed masses and
increase the observed momentum widths of the primary fragments[36]. The
contribution from sequential decay of primary fragments to the heavy
projectile fragments are less than that to the light projectile fragments, the
widths of the transverse momentum distributions of light projectile fragments
are greater than that of the heavy projectile fragments. So the transverse
momentum distribution width deceases with the increase of the charge of the
projectile fragment, which is shown in Figure 7.
Figure 8 shows relation of the averaged transverse momentum per nucleon and
the charge of projectile fragment for interactions of 56Fe and C, Al, and CH2
targets. The averaged transverse momentum per nucleon decreases with the
increase of fragment size, no obviously target size dependence is observed in
present investigation. Figure 9 shows the transverse momentum distribution of
all fragments for interactions of 56Fe and C, Al, and CH2 targets, no
obviously target size dependence is also observed. The heavy fragment comes
from peripheral collisions, that is collisions with larger impact parameter.
The light fragment comes from central and semi-central collisions, that is
collisions with smaller impact parameter. According to the participant-
spectator model[37], with the increase of impact parameter the overlapped
region decreases, the communication between participant and spectator
decrease. This results in the decrease of the excitation energy of projectile
fragments, so the average of transverse momentum of fragment is decreased.
Figure 7: The transverse momentum distributions of fragments with charge Z=6,
18, and 24, respectively. Figure 8: Relation of the averaged transverse
momentum and the charge of projectile fragment for interactions of 56Fe and C,
Al, and CH2 targets. Figure 9: Transverse momentum distribution of all
fragments for interactions of 56Fe and C, Al, and CH2 targets.
Figure 10: The cumulative $p_{t}^{2}$ distribution of projectile fragments
with charge $Z=6$ (a), $Z=18$ (b), $Z=24$ (c), and $Z=25$ (d).
Based on the participant-spectator concept and the fireball model[38], the
large number of swept out nucleons combined with an anticipated fairly large
number of interactions per particle is presumably responsible for the quasi-
equilibrated system, i.e. the fireball which can then be described in term of
mean value and statistical (Maxwell-Boltzmann) distribution. If we assume that
the emission of projectile fragments is Maxwell-Boltmann distribution in the
projectile rest frame with a certain temperature T, then the integral
frequency distribution of the square of the transverse momentum per nucleon is
$lnF(>p_{t}^{2})=-\frac{A}{2M_{p}T}p_{t}^{2}$ (5)
where A is the mass number of fragment, $M_{p}$ is the mass of proton. The
linearity of such a plot would be strong evidence for a single temperature of
emission source.
Figure 10 shows the cumulative plots of $F(>p^{2}_{t})$ as a function of
$p^{2}_{t}$ for projectile fragments from the fragmentation of 56Fe on C, Al
and CH2 targets. All of the plots can be well fitted by a single Rayleigh
distribution of the form
$F(p_{t}^{2})=Cexp(\frac{-p_{t}^{2}}{2\sigma^{2}}),$ (6)
where $\sigma=\sqrt{2/\pi}<p_{t}>$, which is related to the temperature of
fragment emission source, $T=\sigma^{2}A/M_{p}$. The fitting parameters
including $\chi^{2}$/DOF and the temperature of the emission source are
presented in Table 2. Because CR-39 detector can not identify the mass numbers
of projectile fragments, we use the mass number of stable nucleus to calculate
the temperature of projectile fragment emission source. The influence from
isotope is less than $1\%$ when the abundance of isotope is considered. The
dependence of the temperature of emission source on the size of fragment for
different targets is shown in figure 11. From the results of Table 2 and
figure 11 we can conclude that the temperature of projectile fragment emission
source decreases with the increase of the charge of fragment for the same
targets. The temperature increases with the increase of target size for
emission of fragment with change $Z=6$ and 18, but for emission of fragment
with change $Z=24$ and 25 this dependence is not obvious.
Table 2. Values of fitting parameters of $p_{t}^{2}$ distribution using
Rayleigh distribution and the temperature of the emission source. Z Target C
$\sigma^{2}$ ((MeV/c)2) T (MeV) $\chi^{2}$/DOF 25 CH2 $96.52\pm 5.35$
$31.93\pm 1.69$ $1.87\pm 0.10$ 0.831 C $100.16\pm 5.71$ $19.47\pm 1.11$
$1.14\pm 0.07$ 0.235 Al $100.56\pm 16.91$ $14.50\pm 2.78$ $0.85\pm 0.16$ 0.033
24 CH2 $96.77\pm 5.95$ $46.51\pm 3.01$ $2.58\pm 0.17$ 1.398 C $101.72\pm 9.24$
$32.52\pm 2.89$ $1.80\pm 0.16$ 0.575 Al $99.56\pm 17.90$ $27.25\pm 5.00$
$1.51\pm 0.28$ 0.159 18 CH2 $99.45\pm 9.32$ $108.00\pm 9.34$ $4.61\pm 0.40$
0.142 C $95.86\pm 14.60$ $133.30\pm 22.31$ $5.68\pm 0.95$ 0.310 Al $94.08\pm
33.96$ $276.24\pm 127.42$ $11.78\pm 5.43$ 0.108 6 CH2 $100.76\pm 11.64$
$1168.60\pm 156.12$ $14.95\pm 2.00$ 0.119 C $94.57\pm 12.88$ $1541.40\pm
221.88$ $19.72\pm 2.84$ 0.400 Al $102.55\pm 27.41$ $1783.50\pm 523.46$
$22.82\pm 6.70$ 0.065
Figure 11: The temperature derived from the distribution of $p_{T}^{2}$.
The temperature of projectile fragment emission source have been investigated
by different collaborations[30,31,39-44]. ALADIN Collaboration studied the
slope temperature ($T_{slope}$) for spectator decay as a function of the
fragment mass ($A_{frag}$) for spectator decays following 197Au on 197Au
collisions, they found that there is a rapid increase of $T_{slope}$ with
fragment mass which saturates for $A_{frag}\geq 3$ around $T_{slope}\sim 17$
MeV[39, 45]. EOS Collaboration also studied the variation of remnant
temperature with the charged particle multiplicity, they found that the
remnant temperature increases with increase of the charged particle
multiplicity and the maximum is about $15.6\pm 0.47$ MeV[46]. These maximum
temperature is the same as our results within experimental errors for emission
of fragment with charge $Z=6$.
According to the participant-spectator concept, it is assumed that when the
interaction of projectile and target nuclei takes place, the projectile and
target sweep out cylindrical cuts through each other. During the separation of
the spectators from the participants, there is some intercommunication, which
results in the excitation of the spectators. This excitation strongly depends
on the contacted area of the colliding system. The heavier fragment is
corresponding to the large impact parameter and small contacted areas, the
lighter fragment is corresponding to the smaller impact parameter and large
contacted areas. So the excitation energy of the heavier fragment is less than
that of the lighter fragment, which results in the temperature of emission
source of heavier fragment is less than that of the lighter fragment.
## 4 Summary
The emission angular distribution and the transverse momentum distribution of
projectile fragments produced in fragmentation of 56Fe on C, Al, and CH2
targets are studied in present investigation. It is found that for the same
target the average value and width of angular distribution decrease with
increase of the projectile fragment charge, and for the same projectile
fragment the average value of the distribution increases and the width of the
distribution decreases with increasing the target charge number. The
transverse momentum distribution of projectile fragment can be explained by a
single Gaussian distribution and the averaged transverse momentum per nucleon
decreases with the increase of the charge of fragment, and no obvious
dependence of transverse momentum on target size is observed. The cumulated
squared transverse momentum distribution of fragment can be well explained by
a single Rayleigh distribution. The temperature parameter of emission source
of projectile fragments decreases with the increase of the size of projectile
fragments.
## 5 Acknowledgements
This work has been supported by the Chinese National Science Foundation under
Grant Nos: 11075100 and 10975019 and the Natural Foundation of Shanxi Province
under Grant 2011011001-2, the Shanxi Provincial Foundation for Returned
Overseas Chinese Scholars, China (Grant No. 2011-058), the Foundation of
Ministry of Personnel of China for Returned Scholars Grant MOP2006138, and the
Fundamental Research Funds for the Central Universities. We gratefully
acknowledge the staffs of the HIMAC for providing the beam to expose the
stacks.
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|
arxiv-papers
| 2013-09-17T13:03:29 |
2024-09-04T02:49:51.018128
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Y.J. Li, D.H. Zhang, S.W. Yan, L.C. Wang, J.X. Cheng, J.S. Li, S.\n Kodaira, N. Yasuda",
"submitter": "Dong-Hai Zhang",
"url": "https://arxiv.org/abs/1309.4294"
}
|
1309.4307
|
# Two rings but no fellowship: LoTr 1 and its relation to planetary nebulae
possessing barium central stars.
A.A. Tyndall1,2, D. Jones2, H.M.J. Boffin2, B. Miszalski3,4, F. Faedi5, M.
Lloyd1, J.A. López6, S. Martell7, D. Pollacco5, and M. Santander-García8
1Jodrell Bank Centre for Astrophysics, School of Physics and Astronomy,
University of Manchester, M13 9PL, UK
2European Southern Observatory, Alonso de Córdova 3107, Casilla 19001,
Santiago, Chile
3South African Astronomical Observatory, PO Box 9, Observatory 7935, South
Africa
4Southern African Large Telescope. PO Box 9, Observatory 7935, South Africa
5Department of Physics, University of Warwick, CV4 7AL, UK
6Instituto de Astronomía, Universidad Nacional Autónoma de México, Ensenada,
Baja California, C.P. 22800, Mexico
7Australian Astronomical Observatory, North Ryde, 2109 NSW, Australia
8Observatorio Astronómico National, Madrid, and Centro de Astrobiología, CSIC-
INTA, Spain
E-mail: [email protected]
(Accepted xxxx xxxxxxxx xx. Received xxxx xxxxxxxx xx; in original form xxxx
xxxxxxxx xx)
###### Abstract
LoTr 1 is a planetary nebula thought to contain an intermediate-period binary
central star system ( that is, a system with an orbital period, P, between 100
and, say, 1500 days). The system shows the signature of a K-type, rapidly
rotating giant, and most likely constitutes an accretion-induced post-mass
transfer system similar to other PNe such as LoTr 5, WeBo 1 and A70. Such
systems represent rare opportunities to further the investigation into the
formation of barium stars and intermediate period post-AGB systems – a
formation process still far from being understood. Here, we present the first
detailed analyses of both the central star system and the surrounding nebula
of LoTr 1 using a combination of spectra obtained with VLT-FORS2, AAT-UCLES
and NTT-EMMI, as well as SuperWASP photometry.We confirm the binary nature of
the central star of LoTr 1 that consists of a K1 III giant and a hot white
dwarf. The cool giant does not present any sign of s-process enhancement but
is shown to have a rotation period of 6.4 days, which is a possible sign of
mass accretion. LoTr 1 also presents broad double-peaked H$\alpha$ emission
lines, whose origin is still unclear. The nebula of LoTr 1 consists in two
slightly elongated shells, with ages of 17 000 and 35 000 years, respectively,
and with different orientations. As such, LoTr 1 present a very different
nebular morphology than A70 and WeBo 1, which may be an indication of
difference in the mass transfer episodes.
###### keywords:
planetary nebulae: individual: LoTr 1, WeBo 1, A66 70 – stars: AGB and post-
AGB – stars: binaries: general – stars: chemically peculiar
††pagerange: Two rings but no fellowship: LoTr 1 and its relation to planetary
nebulae possessing barium central stars.–References††pubyear: 2013
## 1 Introduction
Figure 1: Deep, narrowband images of PN LoTr 1, in (a) [O iii]$\lambda$5007 Å,
and (b) H$\alpha$+[N ii]$\lambda$6584 Å. North is to the top of the image,
East is left. The central star is visible at $\alpha$ = 05:55:06.6, $\delta$ =
$-$22:54:02.4. Overlaid on image (b) are the slit positions of the spatio-
kinematic data presented in section 2.1.2. Slit 1 (E-W) was acquired using
NTT-EMMI and is centered on the central star. Slits 2–4 (N-S) and 5–8 (E-W,
width of slits 6, 7 and 8 exaggerated to clearly show the positioning due to
overlapping with slit 1) were acquired using AAT-UCLES. All slits were taken
in [O iii].
The interaction between the progenitor star and a binary or planetary
companion is believed to shape the resulting planetary nebula (PN), and in
some cases even thought to be almost essential for a PN to form (Moe & De
Marco, 2006). The shaping influence of a common-envelope evolution has been
studied extensively (see e.g. Tyndall et al., 2012, Jones et al., 2010), and
understood in terms of either a collimated fast wind (CFW) carving out an
axisymmetric nebula (Soker & Rappaport, 2000) or the ejected common-envelope
(CE) forming an equatorial density enhancement (Nordhaus & Blackman, 2006) as
required by the ‘Generalised Interacting Stellar Winds’ model (Kahn & West,
1985). However, very little is known about intermediate-period (P=100-1500
days, van Winckel et al., 2009) post-Asymptotic Giant Branch (AGB) binaries,
including their effect on PNe formation and morphology, due to lack of
observations. The intermediate period binaries fall between post-CE systems
(Miszalski et al., 2009) and visually resolved systems (e.g. Ciardullo et al.,
1999). Soker (1997) claimed that these interacting systems are the most likely
CSPNe (central stars of planetary nebulae) to form the classical ‘butterfly’,
or bipolar, morphologies, but with few systems known and limited
investigations, this has yet to be confirmed. Only by finding and studying
CSPNe with intermediate periods can we substantiate this claim and relate the
processes at work in formation of PNe by both post-CE and intermediate-period
CSPNe.
The planetary nebula (PN) LoTr 1 ($\alpha$ = 05:55:06.6, $\delta$ =
$-$22:54:02.4, J2000) was first discovered by A.J. Longmore and S.B. Tritton
with the UK 1.2-m Schmidt telescope (Longmore & Tritton, 1980). It is
generally noted that LoTr 1 belongs to the so-called ‘Abell 35-type’ group
(Bond et al., 1993) of PNe showing evidence of a binary central star system
consisting of a cool central star (a rapidly rotating subgiant or giant), and
an optically faint hot companion (a white dwarf with effective temperature,
$t_{{\rm eff}}$ $\sim$ 100 kK) since these giant stars are too cool to ionise
the surrounding nebula. Four PNe fell into this category: Abell 35 (hereafter,
A35), LoTr 5 (Thevenin & Jasniewicz, 1997), WeBo 1 (Bond et al., 2003) and
Abell 70 (Miszalski et al., 2012, hereafter, A70). However, Frew (2008)
determined that A35 is most likely not a true PN, but rather a Strömgren zone
in the ambient interstellar medium (ISM). This claim is substantiated by
Ziegler et al. (2012), who find that the central star may in fact have evolved
directly from the Extended Horizontal Branch to the White Dwarf (WD) phase (a
so-called AGB-manqué star). As such, we choose not to consider Abell 35 in our
comparisons among this group.
Another common factor amongst this particular group of PNe is evidence for the
existence of ‘Barium (Ba ii) stars’ (Bidelman & Keenan, 1951) – population i
G/K-type AGB stars that show an over-abundance of carbon and s-process
elements, in particular barium (Thevenin & Jasniewicz, 1997). A now-canonical
model for the formation of these Ba ii stars states that they form not through
CE evolution111We note, however, the existence of some exceptional systems
that also experience similar enrichment in close binaries, i.e. the Necklace
nebula (Miszalski et al., 2013), and which are most likely linked to dwarf
carbon stars., as is the case for close binaries found within PNe, but rather
via a wind-accretion scenario (Boffin & Jorissen, 1988). Here, the future Ba
ii star is polluted whilst on the main sequence by the wind of its companion
(Luck & Bond 1991; the companion having dredged up these s-process elements
during its thermally pulsing AGB phase), but with the system remaining
detached. After the envelope is ejected to form the surrounding nebula, the
AGB star evolves into a WD, while the contaminated star retains its chemical
peculiarities to form the remnant Ba ii star.
One important prediction to come out of the wind-accretion model is that the
accreting star, i.e. the future Ba ii star, also accretes angular momentum
from the companion to become a rapid rotator (Jeffries & Stevens, 1996; Theuns
et al., 1996). Indeed, photometric monitoring of LoTr 5 (Thevenin &
Jasniewicz, 1997) and Webo 1 (Bond et al., 2003) has revealed that their cool
components are in fact rapid rotators with a rotation period of a few days,
thus providing further evidence for this formation scenario. This is further
evidenced by the fact that Montez et al. (2010) found that the x-ray emission
from the binary central stars of LoTr 5 is most likely due to the
chromospheric activity from a spun-up companion.
Most Ba ii stars are not observed to be within PNe, almost certainly because
the lifetime of the PN is very short with respect to the lifetime of the
stellar system. However, this does not completely rule out the possibility
that some Ba ii stars may be formed without passing through a PN phase.
Indeed, there are a few examples of field stars that have been shown to
consist of a rapdidly rotating cool star linked to an optically faint hot
component in a similar fashion to the Abell-35 group, but without current
evidence for a surrounding PN. As mentioned by Bond et al. (2003), HD 128220
(P = 872 days) is one such system, made up of an O subdwarf and a G0 giant
companion. O subdwarfs overlap with CSPN in terms of their log g and $t_{{\rm
eff}}$, implying that they too are found to be in a post-AGB phase of
evolution (Howarth & Heber, 1990). 56 Pegasi is another example, whereby it
possess a system consisting of a K0 giant and a hot white dwarf companion with
evidence for an overabundance of Ba ii, and has an estimated orbital period of
111 days (Griffin, 2006). It is possible that both of these systems _have_
gone through a PN phase in the past, but it has since dissipated into the
surrounding ISM. This fact, as well as the knowledge that in Ba ii central
stars the chemical pollution process happened very recently – either during or
immediately prior to the formation of the PN – makes it highly important to
study such systems as A70 and WeBo 1, as it will allow us to gain greater
insight into both the s-process within AGB stars and mass-transfer mechanisms.
In this paper, we present photometric and spectroscopic observations of LoTr 1
and its central star along with complementary data of A70 and WeBo 1, in order
to try to relate the evolutionary processes of these three systems222PN K 1-6
(Frew et al., 2011) is another potential candidate for this group of objects,
as it also shows evidence of possessing both a G- or K-type giant (inferred
from imagery, looking at both optical and 2MASS near-IR colours) and a very
hot sub dwarf or white dwarf (inferred from GALEX archival images) at its
core. However, no stellar spectroscopy is available as yet to look for signs
of chemical enrichment (Frew et al., 2011). More recently, Miszalski et al.
(in press) presented evidence for a carbon and s-process enriched giant at the
centre of the planetary nebula Hen 2-39. Due to it being a newly investigated
system, Hen 2-39 is not included in this study either., and see if they belong
to a common “fellowship” of s-process enriched cool CSPN inside ring-like
nebulae.
Figure 2: PV arrays showing reduced, longslit [O iii]$\lambda$5007 Å NTT-EMMI
and AAT-UCLES spectra of LoTr 1, both aligned E-W, to show the overall nebular
structure. West is to the top of the array. The velocity axis is heliocentric
velocity, $V_{{\rm hel}}$. The display scale has been modified to highlight
the spatio-kinematic features referred to in the text. Cross-section 0″
defines where the central star is found. The gap between the two CCDs in
figure (a) is visible as a black strip across the frame at cross-section 58″
(see text for details). Figure 3: PV arrays showing reduced [O
iii]$\lambda$5007 Å AAT-UCLES spectra of LoTr 1. The velocity axis is
heliocentric velocity, $V_{{\rm hel}}$. The display scale has been modified to
highlight the spatio-kinematic features referred to in the text. Slits 2–4 are
N-S (North is to the top of the array), slit 5 is E-W (West is to the top of
the array). Cross-section 0″ defines where the central star is found. The
continuum of a field star is visible at cross-section $-$8″ in figures (a) and
(b), and at cross-section +18″ in figure (d).
## 2 Observations and Analysis
### 2.1 LoTr1
#### 2.1.1 Imaging
The deep [O iii]$\lambda$5007 Å and H$\alpha$\+ [N ii]$\lambda$6584 Å images
shown in figures 1(a) and (b) were acquired on 2005 March 03 using the red arm
of the European Southern Observatory (ESO) Multi-Mode Instrument (EMMI; Dekker
et al. 1986), mounted on the 3.6m New Technology Telescope (NTT) of the La
Silla Observatory. EMMI was used with the mosaic of two MIT/LL CCDs, 2048
$\times$ 4096 15 $\mu$m pixels. The exposure time, $t_{{\rm exp}}$, in each
filter was 1800s and the binning set to $2\times 2$ ($\equiv$ 0.33″ pixel-1).
The seeing was $\sim$0.7″.
The images show LoTr 1 to have an apparent double-shell structure, with the
central shell having a circular profile with an angular diameter of
47″$\pm$4″, and the outer shell a more irregular, but still roughly circular,
appearance, with a diameter of 2′26″$\pm$4″. Furthermore, in [O
iii]$\lambda$5007 Å the outer shell appears brighter in the northwest and
southeast – this, coupled with the slight deviation from circular symmetry,
could be considered evidence for an inclined, elongated structure, where the
brighter areas result from a projection effect.
Interaction with the ISM is also frequently invoked to explain deviation from
symmetry and brightening of the nebular shell (e.g., Jones et al. 2010),
however this would not produce the apparent axisymmetry in the nebular
brightening (as the nebula would only interact strongly with the ISM in the
direction of its motion; Wareing et al. 2007).
#### 2.1.2 Nebular spectroscopy
Longslit echelle spectroscopy was carried out using AAT-UCLES and NTT-EMMI,
focusing on [O iii] emission over eight different slit positions (see figure 1
(b)) in order to gain velocity profiles across a good sample of the nebula. It
is important to ascertain the true three-dimensional shape of the surrounding
nebula in order to fully understand and constrain the shaping process, and
subsequently give a clearer insight into mass-transfer mechanisms. Imaging is
insufficient on its own due to a degeneracy between PN inclination and
morphology – for example, when the symmetry axis of a bipolar nebula is
aligned perpendicular to the plane of the sky, it makes the nebula appear
spherical (Kwok, 2010). If the nebular morphology is classified incorrectly
then the fraction of aspherical PNe possessing a binary core will also end up
being inaccurate, and this information is already exceedingly limited for
longer period binaries due to selection biases for those stellar systems
possessing wider orbits. Longslit spectroscopy can be used to acquire
spatially resolved velocity maps of the constituent parts of the nebula in
order to recover the ‘missing’ third dimension of the morphology one cannot
gain from imagery alone. The resulting spectra are then plotted as position-
velocity (PV) arrays.
On 2005 March 03, a spectrum was acquired from the nebula using EMMI in its
single order echelle mode, employing grating #10 and a narrowband [O
iii]$\lambda$5007 Å filter to prevent contamination from overlapping echelle
orders. The maximum slit length of 330″ and a slit width of 1″ was used to
give a resolution, R $\sim$ 54,000 (5.5 km s-1). A 1800s exposure was taken at
a P.A. of 90∘ crossing the central star (slit position 1 in figure 1(b)), and
the data were 2x2 binned to give a spatial scale of 0.33″ and a velocity scale
of 3.9 km s-1 per pixel. The seeing was $\sim$0.7″.
On 2005 January 14, spectra were acquired from the nebula of LoTr 1 using the
79 lines/mm grating on the UCL Coudé Echelle Spectrograph (UCLES) of the
Anglo-Australian Telescope (AAT). UCLES was operated in its longslit mode with
a maximum slit length of 56″ and a slit width of 1.97″ to give a resolution R
$\sim$ 20,000 (15 km s-1). The EEV2 CCD (2048$\times$4096 13.5 $\mu$m pixels)
was used with binning of 2$\times$3, resulting in a pixel scale of 3.88 km s-1
pixel-1 in the spectral direction and 0.48″ pixel-1 in the spatial direction.
1800s exposures were taken at five different slit positions (shown in figure
1(b)) using a narrowband filter to isolate the 45th echelle order containing
[O iii] emission line profile. Slits 2–4 were taken at a position angle (P.A.)
of 0∘, and slit 5 was taken at a P.A. of 90∘. The seeing during the
observations was $\sim$2″. A further three slits (slit positions 6–8 in figure
1(b))were acquired using the same instrument and CCD on 2013 January 3 with a
seeing of $\sim$1.5″ and a binning of 2$\times$2 ($\equiv$ 0.32″ pixel-1 in
the spatial direction). Here, a slit width of 1″ was employed for slits 6 and
8 (R $\sim$ 45,000 $\equiv$ 6.7 km s-1) and 1.5″ for slit 7 (R $\sim$ 30,000
$\equiv$ 10 km s-1)333The slit width for slit 7 was altered to try to match
the seeing conditions during the night. However, as all observations were
carried out with the same binning in the spectral direction and the same
grating (i.e. approximately the same dispersion), they can all be used to
qualitatively assess the spatio-kinematic structure of the nebula (varying on
scales $\gg$ the slit-widths employed, due to the angular size of the nebula
also being $\gg$ the slit-width)..
All the spectra were cleaned of cosmic rays and debiased appropriately. The
UCLES spectra were wavelength calibrated against a ThAr emission-lamp,
rescaled to a linear velocity scale appropriate for the [O iii]$\lambda$5007 Å
emission, and corrected to heliocentric velocity, $V_{{\rm hel}}$. Due to the
optical set-up of EMMI, it was necessary to perform the wavelength
calibrations using a long-exposure (3600s) ThAr emission lamp at the start and
end of the night to gain a good number of arc lines, before cross-correlating
with shorter-exposure (200s) Ne lamps taken immediately after each
observation, to account for any drift due to telescope and instrument flexure
(with small shifts accounted for with a linear correction).
The reduced nebular spectra of LoTr 1 are presented in figures 2 and 3, as
position-velocity (PV) arrays. In each PV array, cross-section 0″ defines
where the central star is found. In figure 2 (a), the bright lines located
around the central star at cross-sections +3″ and $-$3″ are most likely
artefacts due to the comparative difference in brightness between the central
star and the nebula (Jones et al., 2010).
The closed velocity ellipses shown in the PV arrays presented in figures 2
(a), 3 (c), and the central PV array of 2 (b) (i.e. representing slit
positions 1, 4 and 7), have a major axis which has the same length as the
diameter of the inner shell (see section 2.1.1). This indicates that the
nebular structure is indeed an isolated, closed shell rather than a projection
effect related to a bipolar structure being viewed end-on. No significant
asymmetries are observed in these velocity ellipses, which are consistent with
a spherical shell or an elongated ovoid viewed directly along the symmetry
axis. The expansion velocity, $V_{{\rm exp}}$ for the inner shell is measured
to be 17$\pm$4 km s-1 at the location of the central star, while the $V_{{\rm
exp}}$ for the outer shell is measured to be 25$\pm$4 km s-1 falling within
the typical range for a PN (Weinberger 1989). Assuming typical expansion
properties for the nebula (i.e. velocity proportional to distance from the
central star), the latter velocity is then the maximum expansion velocity (for
an elliptical shell viewed pole-on) or the uniform expansion velocity of the
shell (in the case of a sphere). The heliocentric systemic velocity, $V_{{\rm
sys}}$ of this central shell was determined to be 14$\pm$4 km s-1.
In slit 1 (figure 2(a)), emission is clearly detected from outside of the
central shell at cross-sections $\sim$40″ and $\sim-$40″, associated with the
outer shell (see section 2.1.1). Here, the emission from the eastern side
appears blue-shifted with respect to the nebular $V_{{\rm sys}}$, and the west
appears red-shifted; this is indicative of an inclined and extended structure,
e.g. an elliptical nebula, where the approaching “lobe” is tilted slightly to
the east of its receding counterpart. Consideration of the other slits
presented in 3 (b), (c) and (d), confirms this asymmetry in the velocity
profile across the nebula, but indicates that the symmetry axis may lie closer
to the northeast-southwest direction than east-west. However, any deviation
from the line of sight must be rather small given the almost circular
appearance of the shell in the images (figure 1).
Determination of the exact structure and inclination of the nebula would
require a more extensive, higher signal-to-noise (given the outer shells
faintness) dataset, covering more of the physical extent of nebula, and a
detailed spatio-kinematical model such as that presented in Jones et al.,
2012. However, it is clear from both the imaging and spectroscopy presented
here that LoTr 1 shows a double-shelled structure with evidence for an
elliptical and slightly inclined outer shell, and a morphologically similar
inner shell but with a different orientation.
Figure 4: SuperWASP photometry of LoTr 1. Upper panel: Periodogram showing a
clear detection of periodicity at $6.3967\pm 0.0005$ days. Middle panel: The
data points folded on the period obtained from the periodogram. Lower panel:
The data binned into phase bins, overlaid is a sinusoidal fit with a peak-to-
peak amplitude of 0.061 magnitudes. Figure 5: (a) Flux-calibrated FORS2
spectrum of the central star system of LoTr 1. (b) Flux-calibrated IUE spectra
of LoTr 1 (solid line), and of NGC 7293 (dashed line) used to determine the
parameters of the white dwarf. Note also the presence of the Mg II emission at
2800 Å which is a sign of chromospheric activity from the cool companion (e.g.
Jasniewicz et al., 1996; Montez et al., 2010)
#### 2.1.3 Stellar photometry
The field of LoTr 1 has been observed by WASP-S between 2006 May 4 and 2012
February 17 with a total of 21407 photometric points obtained with the two
cameras DAS 226 and 228. WASP-S is a wide field survey camera situated at
SAAO, Sutherland, South Africa, and together with its northerly sister,
SuperWASP-N, is designed to obtain extremely accurate photometry of bright
stars in order to search for transits from exoplanets (see Pollacco et al.
2006, for a more detailed description of the facility, the data reduction and
archive). Data mining tools and the public archive are discussed by Butters et
al. (2010). After reduction through the instrument pipeline (Pollacco et al.,
2006), the time series were examined with a simple Lomb-Scargle filter for
periodic signal detection. From the light curves reproduced in figure 4 taken
from Jones et al. (2011), a period of 6.4 days was derived with a peak-to-peak
amplitude of 0.061 magnitude. This is close to the value of 6.6 days with an
amplitude of 0.1 magnitude in V, initially derived by Bond et al. (1989).
Given the period, this amplitude is probably too high to be considered due to
irradiation effects (De Marco, 2006), so the most likely explanation is that
the periodicity is following spots on a cool star – the signature of a rapid
rotator, analogous to the systems of LoTr 5 (5.9$\pm$0.3d, Thevenin &
Jasniewicz, 1997) and WeBo 1 (4.69$\pm$0.05, Bond et al., 2003).
Figure 6: LoTr 1 FORS2 stellar spectrum (solid line) alongside three synthetic
spectra with [Ba/Fe] = 0.1, [Ba/Fe] = 0.2, and [Ba/Fe] = 0.5, smoothed to
match the FORS2 resolution and plotted over a wavelength range which includes
the Ba ii 4554Å line. Figure 7: H$\alpha$ emission seen in LoTr1.
#### 2.1.4 Stellar spectroscopy
Figure 8: H$\alpha$+[N ii]$\lambda$6584 Å image of PN A70 showing all slit
positions. Slits 1–4 were acquired in [O iii]$\lambda$5007 Å using VLT-UVES,
and slit 5 was acquired in H$\alpha$ and [O iii]$\lambda$5007 Å using SPM-MES.
The central star is visible at $\alpha$ = 20:31:32.2, $\delta=-$07:05:17.0.
On 2012 February 10, a low-resolution spectrum was acquired of the central
star of LoTr 1 using the FOcal Reducer and low dispersion Spectrograph (FORS2;
Appenzeller et al. 1998) with grism 1200g, on the Antu unit (UT1) of the Very
Large Telescope (VLT) based at ESO-Paranal. FORS2 was operated with a maximum
slit length of 6′ and a 0.5″ slit width (R $\sim$ 3000), and at a position
angle of 45°. The data were 2$\times$2 binned to give a spatial scale of 0.25″
pixel-1, with a dispersion on the spectral axis of 0.3Å pixel-1. A single 600s
exposure was taken under seeing of $\sim$ 0.7″.
The stellar data were extracted to 1-D and flux-calibrated against standard
star Hz 4, and are presented in figure 5 (a). All reductions were carried out
using standard starlink routines.
The central star of LoTr 1 was classified as a K1 III-type giant by comparing
the flux-calibrated FORS2 spectrum with UVES POP standard stars (Bagnulo et
al., 2003) that had been rebinned and smoothed to match the resolution of the
FORS2 data. To refine the analysis and derive the possible s-process
overabundance, we used the stellar spectral synthesis code of R. Gray,
spectrum version 2.76444See http://www1.appstate.edu/dept/physics/spectrum/,
with models from Castelli & Kurucz (2003). The best fit parameters were
determined to be $t_{{\rm eff}}$ $=4750\pm$ 150 K and log $g=2.0\pm 0.5$,
which is in agreement with the K1 III spectral type. Using the average value
of the absolute magnitude, Mv = +0.7 for a giant of this spectral type as
given in Allen’s ‘Astrophysical Quantities’ (2000), the flux for the spectrum
shown in figure 5 (a) that allows the derivation of an apparent magnitude, mv
= +12.6 (accounting for 20% slit-losses), and extinction $A_{V}$=0.1285
(Schlegel et al., 1998), one can determine a rough distance to the star of 2.6
kpc and a not-unreasonable radius of the giant of 11.5 R⊙. Alternatively, one
can use the average value of $V$ derived from SuperWASP (see figure 4), i.e.
$V=12.44$, to derive a distance of 2.1 kpc. Note that this observed SuperWASP
magnitude is likely to be contaminated by line emission and close field stars
because of the broad observing band and large (14″) pixels. Therefore, we
think it is preferable to use the distance derived from the flux calibrated
spectra.
To check that the giant and the planetary nebula are indeed linked and not
merely a chance superposition, we computed the radial velocity difference
between the nebular lines in the FORS2 spectrum and the giant star’s
absorption lines, cross-correlating our spectrum (shifted so that the nebular
lines are at zero velocity) with our synthetic K1 III spectrum. The resulting
stellar velocity of 4$\pm$2 km s-1 with respect to the nebula implies that the
cool central star and nebula are physically related.
Figure 9: PV arrays showing reduced VLT-UVES spectra in [O iii]$\lambda$5007 Å
from A70. Figures (a) and (c) show emission from the major axis, figures (b)
and (d) are from the minor axis. Positive spatial offsets are to the northern
(1 and 2) or eastern (3 and 4) ends of the slits.The velocity axis is
heliocentric velocity, $V_{{\rm hel}}$. The display scale has been modified to
highlight the spatio-kinematic features referred to in the text. Cross-section
0″ defines where the central star is found.
The FORS2 spectrum was used to check for signs of Ba ii pollution at 4554Å.
Figure 6 shows the FORS2 spectrum of LoTr 1 plotted alongside three synthetic
spectra with various barium enhancements, from [Ba/Fe] = 0.1 to [Ba/Fe] = 0.5.
An over-abundance of 0.5 or greater classifies the system as a definite Ba ii
star, with a value of 0.2–0.5 being possibly a ‘mild’ Ba ii star (Pilachowski,
1977). Due to the relatively low S/N and resolution of the spectrum, we are
unable to determine a definitive value for the barium abundance in LoTr 1, but
we can clearly state that it is much less than 0.5, and therefore does not
show any measurable barium enhancement. This is contrary to both A70 and WeBo
1 (see Miszalski et al. 2012 and Bond et al. 2002, respectively), which both
possess definite Ba ii stars.
On 2013 January 3, 3$\times$30 minute spectra were acquired from the CSPN of
LoTr 1 using AAT-UCLES with the 79 lines/mm grating, operated in full echelle
mode with a slit width of 1″ to give a resolution R $\sim$ 45,000. All spectra
were reduced using standard starlink routines, corrected to heliocentric
velocity and then summed. Unfortunately, because of the large interorders in
the spectra555UCLES operated in this mode gave non-continuous wavelength
coverage from roughly 5200Å to 8900Å across 19 orders with interorder spacing
of 100–200Å., it is not possible to assess the barium abundance, as none of
the available orders contain the Baii lines at 6141.7Å and 6496.9Å. We have
verified, however, that the iron lines are well fitted with a solar abundance
and our preferred model, in agreement with what we derived from the FORS2
spectrum. Looking at the La ii 6390Å and Yii 6222Å lines, it is also clear
that these elements are not overabundant, confirming the lack of s-process
enhancement in LoTr 1.
The AAT spectra allowed us to find that the stellar component of H$\alpha$ is
in emission (see figure 7). The line is very broad and is clearly double-
peaked, with an equivalent width of about 6Å and the velocity spread is 572
km/s. Using the method of Hodgkin et al. (1995), we derive an H$\alpha$
luminosity $L_{\rm H\alpha}$ = 0.044 L⊙. With the above estimated total
luminosity of the star ($L_{*}\sim$ 60 $L_{\odot}$), this gives a value $\log
L_{\rm H\alpha}/L_{*}=-3.12$.
Such H$\alpha$ double-peaked emission lines have been found in the other stars
we are concerned with here, LoTr 5 (Jasniewicz et al., 1994; Strassmeier et
al., 1997) and A35 (Acker & Jasniewicz, 1990; Jasniewicz et al., 1992), but
the origin is still unknown. Rapidly rotating giants, such as RS CVn or FK Com
stars, are known to have high chromospheric activity which is often revealed
by emission cores in some lines. And indeed, Acker & Jasniewicz (1990) find a
modulation of the H$\alpha$ emission line with the rotation phase, while
Jasniewicz et al. (1992) postulate that the variable double-peaked emission
line is the result of an overlap between an absorption and an emission line at
H$\alpha$, with the possibility for the absorption component to be formed in
the photosphere or through a self-absorption process as in Be stars. However,
the FWHM velocity we measure seems too high to be caused by mass motions
inside the chromosphere, while the luminosity is too small to be due to
accretion. Such double-peaked emission are also sometimes found in symbiotic
stars – detached systems which interact via wind accretion. A few of them, out
of an outburst event, produce bipolar nebulae very similar to planetary
nebulae. It turns out that many symbiotic stars show double H$\alpha$ profiles
(e.g. Schild et al., 1996; Burmeister & Leedjärv, 2009), which might be caused
by a narrow absorption component from the giant overlaid with very broad
H$alpha$ emission from high-velocity jets at the core of the system, or by
disc-like structures. It would thus be of interest to further study the
possible link between LoTr 1 and symbiotic stars.
A spectrum of the central star system of LoTr 1 was acquired using the
International Ultraviolet Explorer (IUE) satellite by Bond et al. (1989).
Presented here in figure 5 (b), it indicates a strong UV continuum of $t_{{\rm
eff}}$ $\geq$ 100 kK. Comparing this IUE spectrum to a known WD within PN NGC
7293 (d = 219${}^{+27}_{-21}$ pc, mv = +13.5; Harris et al., 2007) gave us a
reasonable fit: $t_{{\rm eff}}$ $\geq$ 123 kK, and R = 0.017${\rm\
R}_{\odot}$. Based on values of D = 2′22″$\pm$4″, and $V_{{\rm exp}}$ =
25$\pm$4 km s-1 at a distance of 2.1–2.6 kpc, we derived a kinematical age of
33,000$\pm$9,000 years for the outer nebular shell. For the inner shell, given
values of D = 47$\pm$2″ and $V_{{\rm exp}}$ = 17$\pm$4 km s-1, the age is
derived as 16,000$\pm$6,500 years. Using the white dwarf evolutionary curves
from Bloecker (1995) and assuming an average remnant mass of 0.6${\rm\
M}_{\odot}$, we derived a stellar temperature at the age of the PN of
approximately 120,000 K. The derived radius and age of LoTr 1 are consistent
with Bloecker’s evolutionary curves.
### 2.2 A70
#### 2.2.1 Imaging
The H$\alpha$+[N ii]$\lambda$6584 Å image of A70 shown in figure 8 was
acquired on 2012 July 08, using FORS2 under program ID 0.89.D-0453(A), with an
exposure time of 60s and seeing of 0.8″. At first glance the image shows a
general ring-like appearance similar to that of other ring-like PNe (e.g. SuWt
2: Jones et al. 2010); however, just as noted by Miszalski et al. (2012), a
closer inspection reveals a “ridged” profile more like that of a bipolar
nebula viewed end-on (e.g. Sp 1: Jones et al. 2012). Furthermore, this image
shows in detail the low-ionisation knots first identified by Miszalski et al.
(2012). Many of these structures seem to be akin to the cometary globules seen
in the Helix Nebula (dense condensations of molecular gas embedded in the
ionised nebula, see Meaburn et al. 1992), with knotty heads closest to the
nebula centre and extended tails reaching out towards the outer rim. Extended
material is also visible outside the east-southeasterly edge of the nebular
ring (the emission visible to the north of the ring originates from a
background field galaxy).
Assuming the “ring” is a physical structure rather than a projection effect,
one can deduce the inclination of the nebula by deprojection. The angular size
of the nebula was determined to be 44″$\times$ 38″$\pm$2″, falling in line
with the previously given value of 45.2″x 37.8″ by Tylenda et al. (2003),
giving an inclination of $30$°$\pm 10$°.
Figure 10: PV arrays showing reduced SPM-MES spectra of A70. Emission is
through the horizontal axis (slit 5 in figure 8). The velocity axis is
heliocentric velocity, $V_{{\rm hel}}$. The display scale has been modified to
highlight the spatio-kinematic features referred to in the text. Cross-section
0″ defines where the central star is found.
#### 2.2.2 Nebular Spectroscopy
On 2011 June 10-11, high resolution data of the nebula of A70 were acquired in
[O iii] using grating #3 on the visual-to-red arm of the Ultraviolet and
Visual Echelle Spectrograph (UVES) on Kueyen Unit (UT2) of the VLT (Dekker et
al., 2000), under program ID 087.D-0174(A). UVES was operated in its 30″
longslit mode with a 0.6″ slitwidth (R $\sim$ 70,000, 4.3 km s-1 pixel-1) to
give a spatial scale of 0.17″ pixel-1. A filter was used to isolate the [O
iii] emission lines and prevent contamination from overlapping orders. The
seeing was between 0.5″ and 0.7″ for all observations. Four 1200s exposures
were taken over four different slit positions. Slits 1 and 2 were taken with
VLT-UVES at a P.A. of 160∘ and slits 3 and 4 at a P.A. of 70∘, to line up with
the major and minor axis of the nebula, respectively. Slit 5 was acquired in
both [O iii] and H$\alpha$ on 2011 May 15 using the Manchester Echelle
Spectrograph (MES) mounted on the 2.1-m San Pedro Martir (SPM) telescope based
at the Observatorio Astronomico Nacional in Mexico (Meaburn et al., 2003;
López et al., 2012). The full slit length of 5′ was used with a slitwidth of
150$\mu$m ($\equiv$ 2″, R $\sim$ 30,000), and taken at a P.A. of 90∘. The data
were 2x2 binned to give a spatial scale of 0.75″ pixel-1. The seeing was
$\sim$ 3″
The nebular spectra acquired from VLT-UVES shown in figure 9 show two highly
filamentary components, one red-shifted and one blue-shifted, joined by bright
knots of emission where the slits cross the nebular ring, to form a closed
velocity ellipse in both axes. These filamentary and irregular structures are
typical of disrupted nebulae, where instabilities have begun to structurally
deform the shell (Guerrero & Miranda, 2012), and clearly show that A70 is not
simply an inclined ring but instead has “bubbles” extending in the line of
sight. It is reasonable to assume that these bubbles form a closed (as the
velocity ellipses are closed along both axes) and axisymmetric (the blue- and
red-shifted components are roughly symmetrical) structure. The bright emission
at the extremes of each slit indicate that the nebula may have a cusped waist,
with slits 1 and 2 showing the “crow’s foot”-like structure typical of narrow
waisted nebulae viewed along their symmetry axis (Jones et al., 2012).
However, there is a clear brightening in these regions and those of slits 3
and 4, consistent with a bright ring. We therefore deduce that A70 comprises
of such a bright ring, encircling the waist of a disrupted and faint bipolar
shell.
Using the same spectra, a polar expansion velocity $V_{{\rm exp}}$ for A70 was
calculated to be 39$\pm$10 km s-1. This is in agreement with the value for
expansion velocity of $V_{{\rm exp}}$ = 38 km s-1 given by Meatheringham et
al. (1988), although no error was quoted. A $V_{{\rm sys}}$ of $-$73$\pm$4 km
s-1 was determined for the nebula, which is consistent with the value of
$-$72$\pm$3 km s-1 by Miszalski et al. (2012). The kinematical age of A70 was
determined to be 2700$\pm$950 years kpc-1. Taking the distance to the nebula
to be 5 kpc, as given by Miszalski et al. (2012), this gives an overall
kinematical age for A70 of order 13,400$\pm$4,700 years.
The limited depth and resolution of the MES-SPM spectra offer little extra to
the discussion of the nebular morphology and kinematics, other than to show
that there is no major difference between the emission in H$\alpha$ and that
in [O iii]$\lambda$5007 Å. They are shown here in figure 10 mainly for
comparison with similar data acquired from PN WeBo 1, which will be discussed
in section 2.3.
### 2.3 WeBo 1
Figure 11: Images of PN WeBo 1 in (a) H$\alpha$+[N ii]$\lambda$6584 Å, and (b)
[O iii]$\lambda$5007 Å, with (b) showing the two slit positions. Both slits
were acquired in [N ii] and [O iii] using SPM-MES. The central star is visible
at $\alpha$ = 02:40:13.0, $\delta$ = +61:09:17.0. North is to the top of the
images, East is left.
#### 2.3.1 Imaging
The deep H$\alpha$+[N ii]$\lambda$6584 Å image shown in figure 11 (a) is the
result of coadding 2 $\times$ 120s exposures, each with seeing better than
1.2″, acquired as part of the IPHAS survey (Drew et al., 2005) using the Wide
Field Camera (WFC) on the 2.5m Isaac Newton Telescope based at the
Observatorio Roque de los Muchachos, La Palma. The [O iii]$\lambda$5007 Å
image shown in figure 11 (b) was acquired using the same instrument on 2010
September 9, with an exposure time of 1200s and under seeing of 1.4″.
The images show another ring-like morphology, although structurally different
to A70 (see section 2.2), with a pronounced inner edge and fainter, more
extended emission around its entire circumference – the ring is particularly
diffuse in [O iii]$\lambda$5007 Å, as shown by the lack of a visible inner
edge. Similar extended emission is also found in SuWt 2 (Jones et al., 2010)
and HaTr 10 (Tajitsu et al., 1999), where the ring is actually the waist of an
extended bipolar structure, possibly indicating that WeBo 1 may display the
same morphology but with as yet undetected, very faint lobes. In SuWt 2, Jones
et al. (2010) attribute this extended material to structural and brightness
variations across an irregular toroidal structure; however (particularly in
the light of the H$\alpha$+[N ii] spectra acquired - see figure 12), WeBo 1
shows a much more regular and even ring-like shape, indicating that this is
more likely an intrinsic structural property (i.e. a tear-drop rather than
circular shaped cross-section).
Smith et al. (2007) deprojected the ring of WeBo 1, determining that it is
seen almost edge-on with an inclination of 75°$\pm$3° with an inner-ring
radius of $\sim$ 25″, which is consistent with the dimensions of the ring as
cited by Bond et al. 2003 (64″x 22″), and as measured from the images
presented here (65″$\times$20″$\pm$4″).
#### 2.3.2 Nebular Spectroscopy
On 2010 December 10, spectra were acquired of WeBo 1 in both H$\alpha$+[N ii]
and [O iii]$\lambda$5007 Å using SPM-MES. The maximum slit length of 5′ was
used with a slitwidth of 150$\mu$m ($\equiv$ 2″, R $\sim$ 30000) for each
filter. The data were 2x2 binned to give a spatial scale of 0.75″ pixel-1.
Slit 1 was taken at a P.A. of 353∘ and slit 2 was taken at a P.A. of 263∘ to
cover the major and minor axes of the nebula, respectively. The seeing was
$\sim$1.5″. Due to the H$\alpha$ profiles having high galactic background
emission, only the background-subtracted [N ii] and [O iii]$\lambda$5007 Å are
presented here. The reduced PV arrays are presented in figure 12.
Just as in the imagery (see section 2.3.1), the [N ii]$\lambda$6584 Å profiles
show well defined emission originating from the ring, while the [O
iii]$\lambda$5007 Å is much more diffuse. The [N ii]$\lambda$6584 Å profiles
from both slit positions show strong emission at the inner edge of the ring,
with slightly fainter material (also with a lower velocity dispersion)
reaching out to greater angular extents – consistent with the “tear-drop”
cross-section interpretation presented in section 2.3.1. The [N
ii]$\lambda$6584 Å profile from slit 1 shows little velocity difference
between the emission originating from the two opposing sides of the ring,
indicating that the slit is roughly perpendicular to the P.A. of the nebular
symmetry axis; therefore, the deprojected velocity difference of the two sides
of the ring in slit 2 should offer a good measure of the ring’s expansion
velocity. Taking the nebula inclination to be 75°, the deprojected expansion
velocity for the ring was calculated to be $V_{{\rm exp}}$ = 22.6$\pm$10 km
s-1.
A $V_{{\rm sys}}$ of -6$\pm$4 km s-1 was calculated for the nebula. The
kinematical age of WeBo 1 was determined to be of order 7300$\pm$3700 years
kpc-1. Taking the distance to the nebula to be 1.6 kpc (Bond et al., 2003),
this gives an overall age for WeBo 1 of 11,700$\pm$5,900 years.666Bond et al.
(2003) assume an expansion velocity of 20 km s-1 to derive a similar age of
12,000$\pm$6,000 yrs.
No clear evidence of lobes or extended nebular structure are detected in the
spectra, unlike for A70 (see section 2.2.2); however, faint material detected
inside the ring on the [N ii]$\lambda$6584 Å PV array of slit 1 could be
consistent with such a structure. Deeper, higher resolution spectra are
required to confirm the nature of this emission, but it is safe to say that if
any lobes are present they are significantly fainter than those of A70, as its
lobes were still clearly detected by SPM-MES spectra (see figure 10).
Figure 12: PV arrays showing reduced [N ii]$\lambda$6584 Å and [O iii]$\lambda$5007 Å SPM-MES spectra of WeBo 1. Figures (a) and (b) show the [N ii] and [O iii] emission from Slit 1, (c) and (d) from Slit 2 (see figure 11). North is to the top of the array. The velocity axis is heliocentric velocity, $V_{{\rm hel}}$. The display scale has been modified to highlight the spatio-kinematic features referred to in the text. Cross-section 0″ defines where the central star is found. The continuum of a field star is visible at cross-section +15″ in figures (a) and (b). Table 1: Physical parameters of LoTr 1, A 70 and WeBo 1 Nebula | $v_{exp}$ ( km s-1) | Kinematical age (yrs) | Physical size (pc) | Morphology
---|---|---|---|---
LoTr 1 (inner) | 17$\pm$4 | 17,000$\pm$5,500 | 0.59$\pm$0.05 | spherical/elliptical
LoTr 1 (outer) | 25$\pm$4 | 35,000$\pm$7,000 | 1.86${}^{+0.05}_{-0.09}$ | spherical/elliptical
A 70 | 39$\pm$10 | 13,400$\pm$4,700 | 1.10${}^{+0.06}_{-0.04}$ | Ringed-waist with detected lobes
WeBo 1 | 23$\pm$10 | 11,700$\pm$5,900 | 0.50$\pm$0.03 | Ringed-waist, no lobes detected
## 3 Discussion
### 3.1 The Abell-35 group and PN mimics
Abell 35 - the archetype of the class of PNe discussed in this paper – has
recently been shown to be a PN mimic. It is, therefore, critical to establish
the true nature of the objects considered here before beginning a comparison.
We restrict our analysis to the three PNe presented in this work, and exclude
both Abell 35 (which is no longer considered a true PN) and LoTr 5 (which is a
considerably more complex case and discussed in more detail in Frew, 2008; see
also Graham et al., 2004).
Frew & Parker (2010) present a ‘recipe’ for determining whether we can
classify an object as a true PN or not, parts of which we can apply to the
nebulae presented here:
* •
Presence of a hot, blue central star: All three nebulae presented here show
evidence of excess UV flux that point towards the existence of a hot companion
– see Fig. 5a for LoTr 1, Miszalski et al. 2012 for A70, and Siegel et al.
2012 for WeBo 1.
* •
Nebular morphology: Each PN possesses what we would classify as a ‘typical’ PN
shape, with rings (A70 and WeBo 1) and shells (LoTr 1). A mimic is often more
diffuse.
* •
Systemic velocity: The $V_{{\rm sys}}$ of the central star is consistent with
the $V_{{\rm sys}}$ of the nebula for both LoTr 1 and A70 (see section 2.1.4,
and Miszalski et al., 2012), and so we can say that the observed emission
comes from a true planetary nebula. No $V_{{\rm sys}}$ for the central star of
WeBo 1 has been published.
* •
Nebula expansion: All three nebulae have been shown to have an expansion
velocity typical for a PN (see Table 1).
* •
Nebular diameter: Using the values stated in this paper, the physical
diameters of LoTr 1, A70 and WeBo 1 are all of order $\sim$1 pc (see Table 1):
a sensible value for a PN.
* •
Galactic latitude: Two of the three nebulae are found at high galactic
latitudes of -22° and -25° for LoTr 1 and A70, respectively (WeBo 1 is at and
+1°). PNe are more likely to be found away from the Galactic plane than
isolated Strömgren spheres.
We can thus be confident that the three objects studied in this paper show the
characteristics of bona-fide planetary nebulae, although, as mentioned
earlier, their link with the class of symbiotic stars should also be
investigated further.
### 3.2 Conclusions
From the study conducted in this paper, we have been able to show that LoTr 1
possesses a double-shelled, slightly elliptical morphology of age of
35,000$\pm$7,000 years for the outer shell, and 17,000$\pm$5,500 years for the
inner. We have been able to infer the presence of a K1 III-type giant
($t_{{\rm eff}}$ $\sim$ 4500 K) and hot white dwarf ($t_{{\rm eff}}$ $\sim$
123 kK, R = 0.017${\rm\ R}_{\odot}$) binary system at its core. The cool star
has been shown to be kinematically associated with the nebula, and to have a
rotation period of 6.4 days. Although it was not possible to accurately
determine the [Ba/Fe] value for the central star system, we were able to say
with confidence that LoTr 1 does _not_ show any evidence for an over-abundance
of Ba ii. LoTr 1 also presents double-peaked emission lines, which have been
seen in the other PNe with cool central stars.
Unlike LoTr 1, the PNe A70 and WeBo 1 have both been previously confirmed to
contain a Ba ii-enriched central star system at their core. The two nebulae
are also shown here to display morphologies distinct to that of LoTr 1, with
both possessing ring-like waists and possible extended lobes. The similar
morphologies and chemical enrichment strongly imply that the two have
undergone very similar evolutionary or mass-loss processes. It is possible
that the wind-accretion process involved in the formation of Ba ii stars is
also responsible for the formation of these ring-like morphologies. Although
the CSPN of LoTr 1 does share some common traits with those of A70 and WeBo 1
– namely binarity with a hot- and cool-components, and rapid rotation of the
secondary – both the lack of a significant over-abundance of Ba ii and the
marked difference in nebular morphology would imply a difference in the
evolution of this system. The lack of Ba ii enhancement could be explained by
a difference in progenitor mass, metallicity, or simply quantity of mass
transferred via the same wind-accretion process (the amount of material
accreted is strongly dependent on orbital separation – Boffin & Jorissen,
1988). However, as shown by Boffin & Zacs (1994) only a small amount of matter
is needed to be accreted to make a star appear as a barium star and some mass
must have been transferred as it is required in order to spin up the secondary
to its rapid rotator state. The most obvious explanation, therefore, is that
the mass was transferred at an earlier stage in the evolution of the primary,
i.e. before the thermally-pulsing AGB phase, when the s-process elements are
created and brought to the surface. This would allow us to infer that the AGB
evolution of the primary was cut short by this mass-transfer episode, signs of
which should be detectable in the properties of the WD. We strongly encourage
follow-up observations of the system in order to confirm this hypothesis, and
in particular it would be crucial to determine the orbital period of these
systems. It is, however, important to note that given the inclination of the
LoTr 1 nebula (very close to pole-on), any radial velocity variations of the
central star system would be very difficult to detect, particularly for the
expected period of $\sim$1–3 years, and this therefore requires very high
spectral resolution and stable instruments.
Coming back to the nebula, multiple shells are not uncommon in planetary
nebulae, with 25-60% found to show outer structures (Chu et al., 1987).
However, it is important to distinguish here between ‘halo-like’ shells
(Corradi et al., 2003, 2004), which are extended, generally spherical,
structures attached to the inner shell, and detached outer shells, with the
latter being far less common. This is critical as the haloes are generally
understood to be the ionised remnant of mass lost on the AGB which is now
being swept up to form the inner shell, while the formation mechanism for
multiple, detached shells is still a mystery. Schönberner et al. (1997) show
that this may be possible via a combination of photo-ionisation and wind
interaction, or, alternatively, a binary evolution might be responsible for
rapid changes in mass-loss that could form two distinct shells, such as with
Abell 65 (Huckvale et al., 2013).
LoTr 1 clearly shows a detached outer shell which may have been produced by
rapid changes in mass-loss/transfer in the CSPN. The difference in kinematical
ages between the two shells, of roughly 18,000-20,000 years, is much shorter
than the single star evolutionary timescales on which these changes might
occur.
Perhaps the nature of LoTr 5 should also be investigated further, as it too
does not have an apparent ring-like morphology despite possessing a rapidly
rotating, G5 III-type, Ba ii-rich central star. As mentioned earlier, Montez
et al. (2010) have carried out a study into the x-ray emission emanating from
this system and concluded that it is most likely chromospheric in origin,
implying the presence of a spun-up companion.
The results presented here show that we still have some way to go to fully
constrain mass-transfer mechanisms in intermediate-period binary and post-AGB
systems, and further study of other such similar systems would be highly
beneficial, in particular with regards to Ba ii pollution.
## Acknowledgments
AAT gratefully acknowledges the support of STFC and ESO through her respective
studentships. This work was co-funded under the Marie Curie Actions of the
European Commission (FP7-COFUND). Based on observations made with ESO
telescopes at the La Silla Paranal Observatory under programme IDs:
0.74.D-0373(A), 087.D-0174(A), 087.D-0205(A), 088.D-0750(A), 0.88D-0573(A),
and 0.89D-0453(A). We thank the staff at the ESO La Silla Paranal Observatory
and the Anglo-Australian telescope for their support in the acquisition of the
observations. WeBo 1 images were acquired with the 2.5m Isaac Newton Telescope
located in the Spanish Observatorio del Roque de los Muchachos on La Palma,
Canary Islands, which is operated by the Instituto de Astrofśica de Canarias
(IAC). Both A70 and WeBo 1 analyses are based on observations obtained at the
Observatorio Astronómico Nacional at San Pedro Martir, Baja California,
Mexico, operated by the Instituto de Astronomá, Universidad Nacional Autónoma
de México.
Some of the data presented in this paper were obtained from the Mikulski
Archive for Space Telescopes (MAST). STScI is operated by the Association of
Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555.
Support for MAST for non-HST data is provided by the NASA Office of Space
Science via grant NNX09AF08G and by other grants and contracts.
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|
arxiv-papers
| 2013-09-17T13:48:36 |
2024-09-04T02:49:51.026283
|
{
"license": "Public Domain",
"authors": "A.A. Tyndall, D. Jones, H.M.J. Boffin, B. Miszalski, F. Faedi, M.\n Lloyd, J.A. L\\'opez, S. Martell, D. Pollacco, and M. Santander-Garc\\'ia",
"submitter": "Amy Tyndall",
"url": "https://arxiv.org/abs/1309.4307"
}
|
1309.4501
|
# A fully automatic problem solver with human-style output
M. Ganesalingam and W.T. Gowers
## 1 Introduction and related work
### 1.1 Overview of the paper
The main purpose of this paper is to describe a program that solves elementary
mathematical problems, mostly but not exclusively in metric space theory, and
presents the solutions in a form that is hard to distinguish from solutions
that human mathematicians might write. The following two proofs are examples
of the program’s output.111The program produces LaTeX output, which we
reproduce verbatim here. The first is a proof that if $f:X\to Y$ is a
continuous function and $U$ is an open subset of $Y$, then $f^{-1}(U)$ is an
open subset of $X$, and the second is a proof that if $f:X\to Y$ is an
injection and $A$ and $B$ are subsets of $X$, then $f(A)\cap f(B)\subset
f(A\cap B)$.
> Let $x$ be an element of $f^{-1}(U)$. Then $f(x)\in U$. Therefore, since $U$
> is open, there exists $\eta>0$ such that $u\in U$ whenever $d(f(x),u)<\eta$.
> We would like to find $\delta>0$ s.t. $y\in f^{-1}(U)$ whenever
> $d(x,y)<\delta$. But $y\in f^{-1}(U)$ if and only if $f(y)\in U$. We know
> that $f(y)\in U$ whenever $d(f(x),f(y))<\eta$. Since $f$ is continuous,
> there exists $\theta>0$ such that $d(f(x),f(y))<\eta$ whenever
> $d(x,y)<\theta$. Therefore, setting $\delta=\theta$, we are done.
> Let $x$ be an element of $f(A)\cap f(B)$. Then $x\in f(A)$ and $x\in f(B)$.
> That is, there exists $y\in A$ such that $f(y)=x$ and there exists $z\in B$
> such that $f(z)=x$. Since $f$ is an injection, $f(y)=x$ and $f(z)=x$, we
> have that $y=z$. We would like to find $u\in A\cap B$ s.t. $f(u)=x$. But
> $u\in A\cap B$ if and only if $u\in A$ and $u\in B$. Therefore, setting
> $u=y$, we are done.
The ‘human-style’ output demonstrated by these two examples is the most
visible novel feature of the program, since the problems it solves are all
well within the capacity of existing programs. However, this novelty reflects
a less obvious novelty in the way that the program operates. It would not be
able to produce human-style output if it did not mirror very closely the way
that human mathematicians think. We envisage that the program will be just the
first in a sequence of programs, and we hope that as a result of our detailed
attention to human thought processes, future programs in the sequence will be
able to solve problems that have not previously been solved by fully automatic
provers.
The structure of this paper is as follows. We begin with a brief discussion of
why there has not been more interaction between mathematicians and researchers
in automatic theorem proving. Over the following three subsections, we talk
about human-oriented theorem proving and explain why we see this tradition as
appropriate for what we are trying to do. We then discuss previous work on
natural-language output for theorem-proving programs.
In the second section of the paper, we explain how our program works. We start
by informally presenting an example of the program in action. This is followed
by a description of the ‘moves’ that the program makes. We then present a
second example, this time discussing more explicitly why the program does what
it does. We end the section by explaining how the program writes up its
thoughts to produce the kind of output shown above.
In the third section, we describe an experiment that we conducted in order to
test whether the output of the program was hard to distinguish from genuine
human mathematical writing. We conclude with a brief discussion of our future
plans.
### 1.2 Mathematicians’ attitudes to automated provers
We come to this project as mathematicians who are fascinated by the question
of how people manage to think of proofs of difficult theorems. Since this
question would seem to be fundamental to the whole enterprise of mathematical
research, one might think that large numbers of mathematicians would have
studied it. However, although mathematicians certainly have thought about
their thought processes – the most famous example being Polya polya1957solve ;
polya1954mathematicsi ; polya1954mathematicsii – far less attention is paid
to the topic than one might expect, and papers are routinely written in a
style that appears to do its best to conceal how the ideas they contain were
discovered.
The result is that the work of analysing and understanding in detail how
mathematicians arrive at proofs has to date been mainly carried out by the
automatic theorem proving community. Unfortunately, this work has been largely
ignored by mathematicians: regrets about this have been expressed in many
places in the literature. In a recent article, Bundy discussed this
indifference and suggested several reasons for it bundy2011 . From his reasons
one general theme emerges: the mathematical style of current automated provers
is very different from the style of human mathematics. This is not just a
question of the language such provers use to express problems and their
solutions, but also the nature of the solutions themselves, which are often
very long and use low-level arguments where a human mathematician would use
high-level arguments. Bundy suggests that things are changing (in response to
the last point he suggests using ‘hiproofs’, a hierarchical presentation of
proofs that allows the reader to look at them at several different levels of
granularity) and that it is inevitable that automated provers will in due
course become an indispensable tool for mathematicians. We agree with this,
and one of our motivations is to do what we can to hasten this development.
What could induce mathematicians to be interested in automated provers? In
some cases, the answer is easy: there are mathematics problems that require
large searches or enormously complicated calculations, and computers are
better at these tasks than humans. But most mathematics is not like that, and
for the part that is not, there is a difficulty. Current automated provers are
still a long way from being able to answer questions (except of the specific
kind just mentioned) that might arise in a typical research project, so the
task for the moment is to teach computers to solve much easier questions of a
kind that mathematicians can easily do by hand. This task, though it has
considerable intrinsic interest, has no immediate payoff for mathematical
researchers, so if one wants input from mathematicians who are not specialists
in automated theorem proving, then the barrier to entry will have to be very
low.
For this reason it would be highly desirable to have an automated prover with
the following properties.
1. 1.
(_User-friendly input._) One can input problems without needing to learn a
special-purpose formal language.
2. 2.
(_User-friendly output._) The program will output solutions expressed in the
language that mathematicians customarily use.
3. 3.
(_Informative solutions._) The solutions are not just certificates of truth:
they are arguments of the kind that a good human mathematician might produce.
4. 4.
(_Easy extendability._) One can increase the power of the program by adding
data in the form of facts, problem-solving tips, etc., rather than having to
rewrite the entire program.
In short, it would be wonderful to have an automated theorem prover with which
one could interact in much the way one interacts with a human mathematician.
This is an ambitious goal that will be difficult to achieve, but if it can be
achieved, then it is likely that far more mathematicians will become
interested in automatic theorem proving.
### 1.3 Human-oriented theorem proving: some background
Since the early days of automated theorem proving, there have been two
competing broad approaches. The first, _human-oriented_ , tradition focuses on
analysing human methods and replicating them in programs. The second,
_machine-oriented_ , tradition instead relies on the brute strength of
computers, using extensive search to solve problems in ways that a human would
find difficult to replicate. Given our long-term goal of helping to create a
program that mathematicians will want to use, it is natural that our own work
should belong firmly in the human-oriented tradition. However, this tradition
is distinctly out of favour at present; for example, if you look at the range
of solutions to the standard TPTP test library automated theorem provers tstp1
; tstp2 , you could be forgiven for not realizing that the human-oriented
tradition had ever existed.
Since the human-oriented tradition is less well known these days, we shall
briefly describe it here. Our account is not intended to be comprehensive: we
focus on the strands that most closely relate to our own work. More detailed
accounts of the various periods in theorem proving may be found in
bledsoe1988survey , wos1985overview , mackenzie2001 , pp. 464–469 of
harrison2009book , the introduction to bundy1983 , bundy1999survey , and
cordeschi1996 .
The earliest human-oriented prover was Newell and Simon’s 1956 _Logic Theory
Machine_ newell1957logic , which proved theorems from Russell and Whitehead’s
Principia Mathematica russell1910 , and is often considered to be the first
artificial intelligence program. Its methods were directly modelled on those
of humans: Newell and Simon studied their own behaviour when proving theorems,
and abstracted out _heuristics_ which they then implemented in their program.
Their goals were wider than those of contemporary theorem proving: they hoped
not just to prove theorems but also to ‘[apply] research on complex
information-processing systems … [to] human learning and problem solving’.
Indeed, they used the successor to the Logic Theory Machine, their _General
Problem Solver_ newell1959gps , as the foundation of a book entitled _Human
Problem Solving_ newell1972human .
In the subsequent decade, much of the work on automated theorem proving,
including the General Problem Solver and Gelernter’s _Geometry Theorem Proving
Machine_ , continued in the heuristic, human-oriented tradition, sometimes
drawing directly on mathematicians’ analyses of their own methods in works
such as polya1957solve ; polya1954mathematicsi ; polya1954mathematicsii . This
trend was arrested by Robinson’s highly influential paper _A Machine-Oriented
Logic Based on the Resolution Principle_. Robinson’s _resolution_ was a
single, simple inference rule which turned out to be dramatically more
effective than previously existing methods. Unlike the earlier methods,
resolution had no natural human correlate: formulae needed to be converted
into a not very readable ‘normal form’ before resolution could be applied, and
actual resolution proofs were often difficult or impossible for humans to
follow. However, the hugely increased efficiency of resolution meant that
these were minor concerns.
The next few years saw an explosion of research into resolution. Numerous
improvements were rapidly found, and it seemed at the time that a suitably
advanced version of resolution might solve the problem of automated theorem
proving well enough for all practical purposes. As a result, human-oriented
work was more or less abandoned: while a handful of researchers, including
Newell and Simon, remained focused on human cognition, the majority were
concerned with problem-solving effectiveness and so devoted their attention to
resolution.
According to mackenzie2001 , the optimism of those times was somewhat dimmed
by the discovery of NP-completeness cook1971complexity ; karp1972reducibility
, and the associated realization that it was unlikely that there was a
complete, polynomial time algorithm for theorem proving. Nevertheless, the
majority of researchers continued to focus on machine-oriented methods
throughout the 70s. The notable exception to this was a group run by Woody
Bledsoe at the University of Texas in Austin. Bledsoe’s first steps in the
direction of human-oriented theorem proving may be found in
bledsoe1971splitting , which uses top-level heuristic techniques to split a
problem into smaller subproblems which are then solved by resolution. This
paper was soon followed by a 1972 paper on limit theorems bledsoe1972limit ,
written by Bledsoe and his students Boyer and Henneman. They took an existing
system bledsoe1971splitting and replaced resolution with a new procedure,
which he described as ‘[bearing] a closer resemblance to the proof techniques
of a mathematician than does resolution.’ Among other results, this procedure
was able to prove that differentiability implies continuity and that the limit
of a product of functions is the product of the individual limits. In the
light of the fourth desirable property mentioned above, one feature of the
procedure that is worth highlighting is that it ‘[did] not release its action
unless its need is detected’; that is, it did not interfere with those proofs
in which was not potentially useful. This is in sharp contrast to resolution,
where adding irrelevant axioms can significantly degrade performance. The
paper notes that such heuristics ‘should be sought for other areas of
mathematics’.
Five years later, Bledsoe wrote a seminal paper on mathematical human-oriented
proving, _Non-resolution theorem proving_ bledsoe1977nonresolution . In it, he
described himself as one of a number of researchers who had ‘made the switch’
after finding that resolution and related techniques had encountered enormous
difficulty in solving problems that a human would find easy, and that
substituting human techniques resulted in a program which ‘easily succeeded’.
Bledsoe expressed the opinion that “purely syntactic” methods such as
resolution had reached a plateau, and that further progress would not be made
without using more human techniques. He described the key to progress as
somehow managing to use ‘the knowledge accumulated by humans over the last few
thousand years, to help direct the _search_ for proofs’. The paper went on to
outline the main concepts used in human-oriented proving at the time. A number
of these are very relevant to the work we describe below, including the use of
rewrite rules, forward chaining, typing, a reluctance to expand definitions
unless necessary, and the use of ‘natural’ or ‘goal-oriented’ systems.
Bledsoe continued to work on human-oriented proving for the remainder of his
life; notable achievements include the proof of a number of results using non-
standard analysis bledsoe1977automatic , a proof of the intermediate value
theorem bledsoe1977setvariables , the use of examples to guide proofs
bledsoe1983usingexamples , and analogical proof construction
bledsoe1995analogy . Reading through these papers, one has no sense that
human-oriented proving had reached a plateau: Bledsoe’s approach to automating
the discovery of proofs of mathematical theorems continued to be fertile
throughout his life. Unfortunately, the ‘AI winter’ meant that funding for
mathematical theorem proving dropped sharply soon after Bledsoe ‘made the
switch’. As a result, many of Bledsoe’s students were forced to move to areas
with more direct commercial relevance than mathematics, with the result that
the Bledsoe tradition in mathematical ATP effectively petered out. Bledsoe’s
students did, however, carry his human-oriented tradition into other areas of
theorem proving. Of particular note are Robert Boyer and J. Strother Moore,
who created a family of provers that are used to this day, primarily for
software and hardware verification.
Although many features of the Boyer-Moore provers, NQTHM and ACL2, were driven
by the needs of software verification as opposed to mathematics, two features
of these provers are significant in a mathematical context. First, Boyer and
Moore emphasise the importance of creating a single program that could be used
to tackle a range of problems, whereas Bledsoe constructed a sequence of
loosely related provers, each intended to tackle different problems. Second,
based on intensive inspection of their own proof techniques, Boyer and Moore
introduced an overall architecture, the ‘waterfall’ architecture, which was
much more systematic than the architectures introduced by Bledsoe. The Boyer-
Moore waterfall consisted of four heuristics, ordered in terms of their
relative attractiveness to humans. A Boyer-Moore prover repeatedly applies the
least risky heuristic that can be used, terminating when the result is proved
or when no heuristic can be applied.
The Boyer-Moore provers partially inspired work in the next chapter of human-
oriented theorem proving, the work done by the group run by Alan Bundy at the
University of Edinburgh. Bundy had been influenced by real, human mathematics
from the beginning of his career bundy1975analysing ; bundy1985discovery , and
reports a Bledsoe-like experience of disillusionment with machine oriented
methods bundy1983book , but his major achievements in human-oriented proving
began with a later attempt to reconstruct the behaviour of the Boyer-Moore
prover within a more systematic framework bundy1988proofplans . This led him
to formulate the notion of a _proof plan_ , a computational representation of
a high-level heuristic used by human mathematicians. Because one application
of a proof plan corresponded to many smaller reasoning steps, a theorem prover
that used proof plans could operate with a dramatically reduced search space.
This often led to a proof discovery process that was very similar to that of a
human mathematician. There is an extensive literature on proof plans, which we
will not discuss here, since proof plans are not explicitly used in the work
described below. A description and evaluation of this work may be found in
bundy2000critique .
Meanwhile, on the machine side there was a notable triumph: a result with a
strong claim to be the first solution by an automated theorem prover of an
unsolved problem in mainstream mathematics. The problem was the Robbins
conjecture, that an algebra that obeyed a certain axiom must be a Boolean
algebra, which had been open for six decades when William McCune proved it in
1996 using an automated theorem prover called EQP (for ‘equational prover’)
and eight days of computer time mccune1997solutionof .
This might appear to indicate that the two aims of producing more powerful
provers and of producing provers that will be easy to use by human
mathematicians are in conflict. We do not believe that this is the case, for
reasons that we shall now explain.
### 1.4 Doing more with less
As we have already said, if one wants to write a program with which
mathematicians can interact easily, it is natural to be drawn to the human-
oriented tradition. However, as the history of that tradition makes clear,
that is by no means the only motivation. There is also a conviction amongst
proponents of human-oriented methods that in order to build provers that can
solve certain kinds of problems – roughly speaking, the problems that human
mathematicians typically deal with, apart from the ones that need a lot of
search or heavy calculation – it is _necessary_ to come to grips with human
problem-solving techniques. For example, Bundy says the following in his
influential _Science of Reasoning_.
> Although our science of reasoning might find application in the building of
> high performance automatic theorem provers, the two activities are not co-
> extensive. They differ both in their motivation and their methodology.
>
> I take the conventional motivation of automatic theorem proving to be the
> building of theorem provers which are empirically successful, without any
> necessity to understand why. The methodology is implied by this motivation.
> The theorem prover is applied to a random selection of theorems.
> Unsuccessful search spaces are studied in a shallow way and crude heuristics
> are added which will prune losing branches and prefer winning ones. This
> process is repeated until the law of diminishing returns makes further
> repetitions not worth pursuing. The result is fast progress in the short
> term, but eventual deadlock as different proofs pull the heuristics in
> different directions. This description is something of a caricature. No ATP
> researchers embody it in its pure form, but aspects of it can be found in
> the motivation and methodology of all of us, to a greater or lesser extent.
>
> Automatic theorem provers based on proof plans make slower initial progress.
> Initial proof plans have poor generality, and so few theorems can be proved.
> The motivation of understanding proofs militates against crude, general
> heuristics with low prescriptiveness and no expectancy. The ‘accidental’
> proof of a theorem is interpreted as a fault caused by low prescriptiveness,
> rather then a lucky break. However, there is no eventual deadlock to block
> the indefinite improvement of the theorem prover’s performance. If two or
> more proof plans fit a theorem then either they represent legitimate
> alternatives both of which deserve attempting or they point to a lack of
> prescriptiveness in the preconditions which further proof analysis should
> correct.
>
> […]
>
> Thus, we expect a science of reasoning will help us build better automatic
> theorem proving programs in the long term, although probably not in the
> short term.
Although Bundy’s ‘science of reasoning’ was (at least initially) tightly
focused on the notion of a proof plan, we believe that many of his comments
apply more generally to the advantages of human oriented proving over machine
oriented proving, and especially his conclusion: that in the long term, paying
close attention to human methods will pay dividends.
There are two main reasons for this. A simple one is that a significant
problem with automated theorem provers – some would call it _the_ main problem
– is the danger of combinatorial explosion. It is in the nature of solving a
complex mathematics problem that one throws up other problems that need
solving, which in turn throw up further problems, and so on. If a generous
amount of search is permitted, then this recursive nature of problem solving
naturally leads to the search being iterated, and thus to a combinatorial
explosion. Somehow, in ways that we do not fully understand, humans manage to
avoid this difficulty by keeping search strictly under control. Therefore, it
seems highly plausible that it will not be possible to design provers that can
produce complicated hierarchical proofs without a deep understanding of how
humans manage to do so.
The second reason is more subtle, and best illustrated by means of an example.
Consider the following problem.
###### Problem 1
Let $G$ be a group and let $S$ be a subset of $G$ with the property that if
$x\in S$ and $y\in S$, then $xy^{-1}\in S$. Prove that $S$ is closed under
taking inverses.
This was used as a test problem for a system devised by Reiter in the 1970s
reiter1976models . His system was designed to use models to prune search
trees: if a statement is false in a simple model, then there is no need to
waste time investigating whether it can be proved. It is clear that this
reflects an extremely important aspect of human problem solving: before we
invest time in proving a statement, we like to feel that that statement has at
least some plausibility. Therefore, it is natural to want to incorporate some
kind of checking-against-models facility into a theorem proving program.
Here, translated for human consumption, is what Reiter’s system does. It
begins by taking an element $b\in S$ (if $S$ is empty then the problem is
trivial) and sets itself the goal of proving that $b^{-1}\in S$. The one
property it has available for use is that $S$ is closed under the operation
$(x,y)\mapsto xy^{-1}$, so it uses backwards reasoning to replace the goal by
that of finding $x,y\in S$ such that $xy^{-1}=b^{-1}$.
At this point, it searches for suitable candidates for $x$ and $y$. It has
three statements that these two variables need to satisfy: $x\in S$, $y\in S$
and $xy^{-1}=b^{-1}$. It searches for a solution in a somewhat mechanical way,
ordering the conditions that need to be satisfied and then attempting to
satisfy them one by one. It begins with the order in which the conditions are
presented, so it tries to satisfy the first condition. Since all it knows so
far is that $b\in S$, it takes $x=b$. The first statement is now satisfied, so
it passes to the second and takes $y=b$ for similar reasons. It is now left
needing to prove that $bb^{-1}=b^{-1}$.
Instead of wasting time looking for a proof of this false (in general)
statement, it checks it against a model such as the Klein four-group, rapidly
finding a counterexample. So at this point it backtracks, reorders the
statements that need to be satisfied, and ends up attempting to satisfy what
used to be the third statement, $xy^{-1}=b^{-1}$, without, for the time being,
worrying about whether $x\in S$ and $y\in S$. It finds the simplest solution
to this, namely $x=e$ and $y=b$. Now it is left needing to prove that $e\in
S$, which is straightforward: it needs $u,v\in S$ such that $uv^{-1}=e$ and
the first thing it tries, namely $u=v=b$, works.
Thus, Reiter’s system successfully solves the problem, and even does so quite
efficiently. However, there are at least two ways in which what it does is
strikingly unlike what a human mathematician would do. The first is its
mechanical approach to satisfying the conditions on $x$ and $y$. In general,
finding an object that satisfies two or more properties is a tricky problem,
since it is difficult to focus on more than one condition at a time. Humans
will often begin by assessing which is the hardest condition to satisfy, since
in practice it often happens that the most natural object that satisfies the
hardest condition either already satisfies the easier conditions or can be
easily modified to do so. By contrast, Reiter’s program starts by attempting
to satisfy the conditions in the order in which they are written down, which
leads it to make a wrong guess that no human mathematician would be likely to
make.
If we are just interested in the performance of one program, this is not much
of a problem, since there are only three conditions, so the amount of
backtracking it has to do is very slight. However, if we are content with this
in the short term, then we miss out on search-reducing techniques that may be
useful in the longer term.
The second noticeably non-human move made by the computer is to test the
statement $bb^{-1}=b^{-1}$ against a model (in Bledsoe’s account of this work,
he suggests using the Klein 4-group). To a typical human mathematician, this
statement is just _obviously_ false in general. Of course, it is not much help
to a theorem proving program to be told that a statement is ‘just obvious’, so
let us look more closely at why it is obvious. The answer is that we see very
easily that, by the cancellation law, if $bb^{-1}=b^{-1}$ then $b=e$. And even
at that point we say to ourselves something like, “Not all groups are
trivial,” rather than “The cyclic group of order 2 is a counterexample.” Thus,
if we are content with Reiter’s method, then we miss out on two further
search-reducing strategies: to simplify statements before deciding whether
they are likely to be true, and to make the working assumption that a
sufficiently simple general statement that is not obviously true is almost
certainly false.
That is not all that can be extracted from a close look at how a human would
tackle this problem. Reiter’s system always reasons backwards, but a human
might well look at the hypothesis that $S$ is closed under the operation
$(x,y)\mapsto xy^{-1}$ and reason that since the only element we know to be in
$S$ is $b$, there is no hope of solving the problem without making the
deduction that $bb^{-1}$ is an element of $S$, and therefore that $S$ contains
the identity $e$. And once one has made that deduction, finding $x$ and $y$ in
$S$ such that $xy^{-1}\in S$ is significantly easier than it was before. So we
have another potentially useful technique: if a statement appears to be
essential and can be applied in only one way, then there is no harm in
applying it, even if you cannot see what good it will do.
These remarks are not intended as a criticism of Reiter’s program, which was
successful at what it was trying to do. But they illustrate the general point
that the closer the attention one pays to how humans solve problems, the more
useful information one can extract. This gives us another reason for trying to
make our programs as human as possible. Indeed, we would describe our work as
belonging not just to the human-oriented tradition, but as belonging to the
extreme human end of the machine-human spectrum. In practice, this means
deliberately not allowing ourselves to exploit the speed of computers, for
example by letting them carry out large searches or perform very complicated
calculations. We hope that by submitting to this restriction, we will force
ourselves to develop a number of useful and important techniques while the
problems we are tackling are still relatively simple.
### 1.5 Two kinds of proof
The arguments of the previous section are intended to show that there are
potential advantages to taking a strict human-oriented approach to automated
theorem proving. However, there is no denying that there are also advantages
to using raw computer power. And the fact that it is the machine-oriented
approach that has solved a decades-old unsolved research problem suggests that
its advantages outweigh those of the human-oriented approach. Perhaps the
situation with theorem proving is like the situation with chess, where the
brute force of the computer beats the advanced pattern-recognition skills of
the human grandmaster.
Before we jump to that conclusion, however, we should consider another
possibility: that machine-oriented methods are better for finding certain
kinds of proofs, and human-oriented methods are better for finding other
kinds. The evidence suggests that this is the case. As Kerber puts it
kerber1998 , when talking about machine-oriented programs, “The strength of
these systems is truly remarkable and even open mathematical problems, like
recently the question whether Robbins algebras are Boolean algebras, can be
solved with the assistance of an automated theorem prover. On the other hand,
observing the blind search behaviour of such a system as it fails to solve a
problem that seems trivial to us as humans can be disappointing.”
This phenomenon is not too surprising. Some proofs seem to consist of a
succession of somewhat arbitrary and unpredictable steps, while others can be
discovered by means of what mathematicians would describe as ‘key ideas’. The
proof of the Robbins conjecture comes into the first category: one can think
of it as a combinatorial object that belongs to a vast space – the space of
sequences of well-formed statements that satisfy certain simple rules – inside
which we are looking for a sequence that ends in a certain way. There are no
obvious measures of progress that tell us that some initial segments of
sequences are ‘obviously right’ or at least ‘getting warmer’ (though it is
conceivable that somebody might one day find a more conceptual argument), and
in the absence of such clues there is not much for it but to undertake a huge
search. It is therefore only to be expected that machine-oriented methods have
outperformed human-oriented methods for that problem and for others of a
similar kind.
However, for the majority of proofs that mathematicians find, there is some
kind of ‘story’ to tell of the ideas that give rise to the proof. Typically,
such a story will be a high-level overview of the main difficulty and how it
is overcome, where ‘overcome’ means that the problem is reduced to one or more
problems where that difficulty no longer occurs. Often this reduction is
achieved by means of a well-chosen intermediate statement that turns out to
follow from the initial assumptions and imply the conclusion. The intermediate
statement itself is typically found not by means of a brute-force search but
by a process of approximation: one might make a guess, find that it is
unhelpful, understand why it is unhelpful, and use that understanding to guide
the search for a better intermediate statement. These characteristically human
techniques enable mathematicians to penetrate deep into ‘proof space’, but the
set of proofs that can be discovered in this way forms a tiny fraction of that
space. It seems almost a truism that human methods will be useful for programs
that want to find these special proofs that human mathematicians are so
mysteriously good at finding.
### 1.6 Systems with natural-language output
So far, we have discussed the desirability of a program that will reason in as
similar a way as possible to the way that humans reason. Another desirable
property that we mentioned earlier is that the input and output of the program
should be in the language that mathematicians already use. Several systems
have been developed that use natural language to a greater or lesser extent.
An early example with some similarity to ours is that of Felty and Miller from
1987 feltymiller1987 . They start with a proof tree and convert it into a more
readable form. Their system can also make significant changes to how a proof
is presented. The following is an example of output from their system: it is a
proof that there are infinitely many primes. The function $f$ mentioned in the
proof can be taken to be the function defined by the formula $f(n)=n!+1$: then
the beginning of the proof is asserting some properties of this function that
are subsequently used in the proof (so any other function with those
properties would do just as well).
> Assume $\forall x(f(x)>x)\wedge\forall x\forall
> y(\mathrm{div}(x,f(y))\supset(x>y))\wedge\forall
> x(\neg\mathrm{prime}(x)\supset\exists
> y(\mathrm{prime}(y)\wedge\mathrm{div}(y,x))$. We have two cases. Case 1:
> Assume $\neg\mathrm{prime}(f(a))$. By modus ponens, we have $\exists
> y(\mathrm{prime}(y)\wedge\mathrm{div}(y,f(a)))$. Choose $b$ such that
> $\mathrm{prime}(b)\wedge\mathrm{div}(b,f(a))$. By modus ponens, we have
> $(b>a)$. Hence, $(b>a)\wedge\mathrm{prime}(b)$. Thus, $\exists
> x((x>a)\wedge\mathrm{prime}(x))$. Case 2: Assume $\mathrm{prime}(f(a))$.
> Hence, $(f(a)>a)\wedge\mathrm{prime}(f(a))$. Thus, $\exists
> x((x>a)\wedge\mathrm{prime}(x))$. Thus, in either case, we have $\exists
> x((x>a)\wedge\mathrm{prime}(x))$. Since $a$ was arbitrary, we have $\forall
> n(\exists x((x>n)\wedge\mathrm{prime}(x)))$.
They describe their mechanism for converting the original tree-structured
deductions into readable natural-language text as very simple. It is clear
that with some small changes they could have improved the readability. For
example, they could have replaced $\mathrm{prime}(x)$ by ‘$x$ is prime’,
$\mathrm{div}(x,y)$ by $x|y$ and the symbols for connectives by English words.
However, the result would still have had some slightly odd characteristics –
for instance, no human mathematician would bother to write ‘by modus ponens’ –
that would have betrayed its mechanical origins.
Another program that produced readable text was written by Holland-Minkley,
Barzilay and Constable in 1999 hollandminkley1999verbalization . Their aim was
to create natural-language output from the Nuprl system. This is an
interactive system based on tactics – that is, high-level inference steps –
which is designed to mimic human reasoning. The output from the Nuprl system
is not at all easy for the untrained mathematician to read. However, they
could convert it into language that was considerably closer to what a human
mathematician might write, as the following sample demonstrates.
> Theorem: For integers $a$ and $b$ and natural number $c$, $(a--
> b)--c=a--(b+c)$.
>
> Consider that $a$ and $b$ are integers and $c$ is a natural number. Now, the
> original expression can be transformed to
> $\mathrm{imax}(\mathrm{imax}(a-b;0)-c;0)=\mathrm{imax}(a-(b+c);0)$. From the
> add com lemma, we conclude
> $\mathrm{imax}(-c+\mathrm{imax}(a+-b;0);0)=\mathrm{imax}(a+-b+-c;0)$. From
> the imax assoc lemma, the goal becomes
> $\mathrm{imax}(\mathrm{imax}((a+-b)+-c;0+-c);0)=\mathrm{imax}(a+-b+-c;0)$.
> There are 2 possible cases. The case $0+-c\leq 0$ is trivial. Consider
> $0<0+-c$. Now, the original expression can be transformed to
> $\mathrm{imax}((a+-b)+-c;0+-c)=\mathrm{imax}(a+-b+-c;0)$. Equivalently, the
> original expression can be rewritten as
> $\mathrm{imax}((a+-b)+-c)=\mathrm{imax}(a+-b+-c;0)$. This proves the
> theorem.
In many ways this looks like the kind of continuous prose that a mathematician
would write, though as with Felty and Miller’s system there are a number of
telltale signs of the mechanical origins of the text. For instance, the first
sentence is not quite grammatical: a human would write, ‘Let $a$ and $b$ be
integers and let $c$ be a natural number.’ There is also the trivial point
that mathematicians would write ‘max’ rather than ‘imax’ (trivial because it
would be very easy to change this). There is also a repetitive quality to the
prose that gives it an automatically generated feel: for instance, two
sentences open with ‘Now, the original expression can be transformed to’.
Another system that is often promoted for the readability of its output is
MIZAR. Although this is not an automatic theorem prover, we show an example of
its output for the sake of comparison. The following is a formalized version
in MIZAR of Euclid’s proof that there are infinitely many primes. [From
http://www4.in.tum.de/~wenzelm/papers/romantic.pdf .]
reserve n,p for Nat;
theorem Euclid: ex p st p is prime & p > n
proof
set k = n! + 1;
n! > 0 by NEWTON:23;
then n! >= 0 + 1 by NAT1:38;
then k >= 1 + 1 by REAL1:55;
then consider p such that
A1: p is prime & p divides k by INT2:48;
A2: p <> 0 & p > 1 by A1,INT2:def 5;
take p;
thus p is prime by A1;
assume p <= n;
then p divides n! by A2,NATLAT:16;
then p divides 1 by A1,NAT1:57;
hence contradiction by A2,NAT1:54;
end;
theorem p: p is prime is infinite
from Unbounded(Euclid);
Although this is quite clearly written in a formal language, it is natural
enough that with a bit of effort, one can follow the steps of the argument, at
least if one knows the argument already. And if readability were the only
concern, then once again there would be simple ways of improving the output.
For example, one could remove all the ‘by X’ justifications at the ends of
lines. Of course, readability is _not_ the only concern when a proof is
written in MIZAR. More generally, there are major differences between our aims
and the aims of those who formalize mathematical proofs. Our aim is to write
fully automatic programs with output that meets the typical standards of
soundness of human mathematicians, whereas the proof-verification community
aims to use a great deal of human interaction to produce completely formalized
proofs.
As we have already mentioned, we would ultimately like to contribute to the
creation of a program that mathematicians can use with little or no effort. It
therefore makes sense to set very high standards for the readability of our
output. Another reason for doing this is that we would like to be confident
that our programs really are thinking in a human way. A natural test of this
is that it should be able to produce output that is similar to what a human
produces.
However, passing that test is not by itself sufficient evidence, since in
principle one might be able to create a program that discovers human-style
proofs with the help of a highly mechanical and non-human discovery process.
Such a program could then hide its thought processes and produce nice readable
write-ups. We therefore want to add an extra condition: we would like our
programs to create human-style write-ups _while being faithful to its thought
processes_. By this we mean that the steps that the program takes in order to
discover the proofs it discovers should be directly translated into human
language and should be presented in the order in which they are made. For
example, if a program discovers a short argument after a long search, then the
long search should be presented and not just the short argument.
There is an obvious objection to this, which is that not even human
mathematicians do it. For example, often they will take all sorts of wrong
turns when searching for proofs, and these will not be recorded in the write-
ups. And often they will present, as if out of nowhere, an object that just
happens to make a proof work, hiding from view the process by which they
thought of it.
This is an important point, but the proofs we will discuss in this paper are
not sufficiently difficult to find for the issue to arise in a significant
way. For the proofs the program finds, it is possible to produce output that
is faithful to the program’s thought processes and also quite close to what a
human would write as the final product. When we come to tackle more complex
problems, we plan to have two styles of output. One will be a proof as a human
might write it, and the other will be more like an account, again as a human
might write it, of the proof discovery process. We believe that both styles
are potentially useful to mathematicians.
It is one thing to claim that a program produces output that is similar to
what a human might write; it is another to provide evidence to back up that
claim. Later in the paper, we shall describe an experiment that we carried out
for this purpose.
### 1.7 Soundness
A further motivation for producing human-style write-up is that it gives us a
mechanism for checking soundness: if the program produces proofs that a human
mathematician finds acceptable, then it is doing what we ask of human
mathematicians.
This attitude to soundness is in sharp contrast with most research in
automated theorem proving, especially of the interactive variety. There what
is sought is a _guarantee_ of soundness, which is obtained by carefully
building up a corpus of statements from a small set of axioms, using only
basic rules of logic and statements that have already been established. There
are very good reasons for this: for example, such systems can handle extremely
complicated proofs that are very difficult for humans to check by hand, a
notable recent example being Gonthier’s machine-checked proof of Thompson’s
odd-order theorem gonthier2013machine . They can also check the soundness of
computer hardware, where an absolute guarantee is essential.
However, for a program that produces proofs of relatively easy theorems,
considerations of this kind do not apply. Therefore, we take the attitude to
soundness that a typical human mathematician takes: we want it to be clear to
a human mathematician that the steps the program takes are sound, and we want
the results that the program assumes to be ones that are well established and
in some sense ‘prior’ to the results being proved. Although this does not
provide a formal guarantee of soundness, there are no examples of well
established mathematical results of the level of simplicity we are dealing
with that have been found to be fundamentally wrong, so for all practical
purposes it is a sufficient assurance of soundness.
### 1.8 The other criteria: user-friendly input and easy extendability
The remaining two desirable criteria that a mathematician-friendly program
should have were that it should be possible to input problems without having
to learn a formal language, and that it should be possible to improve the
performance of the program by adding data.
It is much harder to write a program that accepts natural-language input than
it is to write a program that produces natural-language output, since for any
given problem there will be many ways of writing it, and the program has to be
able to handle all of them – or at least, enough of them to make the program
not too restrictive. For the time being, our input takes a form that is not
written in natural language, but is fairly straightforward to write. (For
example, to request that the program show that the pre-image of an open set
$U$ under a continuous function $f$ is open, we supply hypotheses
continuous(f) and open(U) and a target statement open(preimage(f,U)). A more
complicated example is forall x epsilon.(in(x,B) => exists y.(in(y,A) &
lessthan(d(x,y), epsilon))).) However, Barnet-Lamb and Ganesalingam have
recently written a program that is capable of processing many different
statements. Roughly speaking, any mathematical statement that is written at
the level of formality that is typical of the more formal parts of
mathematical textbooks has a good chance of being understood. So it may be
that future versions of the program will be able to accept natural-language
input as well.
As for the possibility of making the program more powerful by adding more data
to it, we are optimistic about this, but it would be premature to make any
claims in this direction. In general, adding mathematical information to a
theorem-proving program can be problematic, because it gives the program more
choice. If the program makes significant use of search, then more choice can
make the searches bigger, and in practice this has meant that some programs
have performed _worse_ when they have more information to go on. The program
has not had this problem, but its library is so small that we cannot at this
stage claim that it never will. The reason we are optimistic is that our
methodology, if we apply it sufficiently rigorously, should ensure that the
problem does not arise. If a program threatens to get bogged down searching
through all sorts of irrelevant results, then that will be a sign that we have
not understood how human mathematicians select relevant results from the large
amount that they know. Understanding that may be very challenging, but if we
cannot fully understand it, then we will try to reach a partial understanding
and make use of that. In other words, we will try to find sufficient
conditions for a use of the program’s data to be human-like, and allow uses
that satisfy those conditions. If the program only ever uses its data in a
human way, then adding data should not degrade its performance, since it does
not degrade the performance of human mathematicians.
## 2 The program
### 2.1 Solving routine problems in a fully human way
The ‘extreme human’ approach we are trying to adopt can be summarized as
follows: _we do not allow our programs to do anything that a good human
mathematician wouldn’t do_. A serious difficulty in trying to design a program
that satisfies this constraint is that while humans do their very best to
avoid search and backtracking, they undoubtedly do at least _some_ search and
backtracking when they are trying to solve difficult problems. What is the
distinction between the kind of search that humans do and the kind that humans
would never do? Indeed, _is_ there a distinction, or is it just a matter of
degree?
These are difficult questions, and rather than try to answer them immediately,
it seems more sensible to try to isolate them by concentrating first on other
difficulties and facing up to this one only when we are forced to do so.
Accordingly, we define a _routine_ problem to be one that a good human
mathematician will typically solve easily without backtracking. If a program
is to satisfy the main constraint, then it too will have to solve routine
problems without backtracking, so we can simply ban it. Although this is a
significant scaling down of ambition, it also makes the project far more
realistic in the short term. Furthermore, the class of routine problems is
large and diverse enough that the challenges it raises are still very
interesting.
### 2.2 Example: closed subsets of complete metric spaces
Before we describe our program, let us look at a routine problem and examine
how a human mathematician would typically solve it. The problem is to show
that a closed subset of a complete metric space is complete. We can state it
more formally as follows.
###### Problem 2
Let $X$ be a complete metric space and let $A$ be a closed subset of $X$.
Prove that $A$ is complete.
The proof discovery process would usually be something like this.
1. 1.
_[Clarify what needs to be proved.]_ We must show that every Cauchy sequence
in $A$ converges in $A$.
2. 2.
_[We must show something about _every_ Cauchy sequence, so pick an arbitrary
one.]_ Let $(a_{n})$ be a Cauchy sequence in $A$.
3. 3.
_[Clarify what now needs to be proved.]_ We are trying to show that $(a_{n})$
converges in $A$.
4. 4.
_[See what we can say about the sequence $(a_{n})$.]_ The sequence $(a_{n})$
is a Cauchy sequence in the space $X$, and $X$ is complete; therefore
$(a_{n})$ converges in $X$.
5. 5.
_[Give a name to the object that we have just implicitly been presented
with.]_ Let $x$ be the limit of the sequence $(a_{n})$.
6. 6.
_[See what we can say about $x$.]_ But $A$ is closed under taking limits, so
$x\in A$.
7. 7.
_[Recognise that the problem is solved.]_ Thus, $(a_{n})$ converges in $A$, as
we wanted.
Our program is designed to imitate these typical human moves as closely as
possible. The result is that what it does can be straightforwardly translated
into an account of its thought processes that could pass for an account of
human thought processes such as the ones laid out above (without the
accompanying commentary).
The architecture in which moves are applied is very similar to that of a LCF-
style interactive theorem prover gordon1979lcf . In simple cases, the program
state consists of a list of statements that can be assumed, which we call
_hypotheses_ , and a list of statements to be proved, which we call _targets_.
Each statement is basically a formula of the (many-sorted) first-order
predicate calculus, but it can also carry annotations that record information
that has been accumulated by the program; for example, a statement may be
tagged to indicate that it has already being used during the course of the
proof. Although information of this kind is logically unnecessary, it is
indispensable in human reasoning and therefore plays an important role in our
automated theorem prover.
The representation diverges from that of a LCF-style prover in more
complicated cases, most notably when a conjunctive target needs to be deduced
from the hypotheses. In this case a LCF-style prover would use a tactic such
as replacing the original goal by two new goals, each of which lists the
hypotheses separately. By contrast, our system represents and displays the
hypotheses only once. This is not done for reasons of efficiency, but rather
to mimic the natural human representation. A human would think in terms of one
ambient primary collection of ‘facts that are known’ (the hypotheses) and a
number of targets that need to be deduced from these; the program reflects
this. In some cases, reasoning will introduce a fact that may be used in the
proof of only one of the targets (and the program represents this
appropriately), but this is very much the marked (i.e. non-default) case and
should not affect the representation of the majority of ‘universally
available’ hypotheses.
An individual _move_ is an operation that transforms a specific problem state
into another state in a sound fashion; thus individual moves correspond to
application of tactics to a specific LCF-style prover state. However, we
insist that our program’s moves faithfully model cognitive processes of human
mathematicians. Thus, many sound transformations of the problem state are not
acceptable as moves, because they correspond to operations that humans would
never perform.
The closest correlate we have to a tactic (as opposed to the application of a
tactic to a specific state) is a _generator_ for a _move type_. An example of
a _move type_ would be ‘forwards reasoning’. Each move type has an associated
_generator_ which, given a state, returns a list of moves from that state. So
for example the generator for the ‘forwards reasoning’ move type accepts a
state and returns a list of all the ‘forwards reasoning’ moves that can be
made from that state. Note that, unlike tactics, generators may return a whole
list of moves.
Because the prover is fully automated, the list of move types is fixed. Move
types are ranked in order of attractiveness or _priority_ , and the basic
operation of the program consists of repeatedly choosing the most attractive
move type that can be applied, generating the moves of that type, and applying
the most attractive one. In practice with the problems we are considering, it
is rare for the most attractive move type to generate more than one move,
except in cases where the moves generated are related by some simple symmetry.
Note that this architecture closely resembles the Boyer-Moore ‘waterfall’
architecture boyermooreacl discussed in §1.3. The waterfall consisted of four
heuristics, ranked in order of their relative attractiveness to humans. A
Boyer-Moore prover repeatedly applies the most attractive heuristic that can
be used, terminating when the result is proved or when no heuristic can be
applied. This directly corresponds to the way in which our prover repeatedly
chooses the most attractive move type possible, terminating when the result is
proved or when no no moves can be made.
We shall now describe what our program does when it tackles the example
problem given above, and then explain in more detail what the move types are,
and what the order of priority is. Like Boyer and Moore, we chose the move
types and their priority by examining our own reactions to many different
problems.
The initial problem state for this problem is as follows.
H1. $X$ is a complete space
H2. $A$ is closed in $X$
T1. $A$ is complete
The first thing the program does is expand the definition of “$A$ is
complete”.
H1. $X$ is a complete space
H2. $A$ is closed in $X$
T2. $\forall(a_{n})$ $(a_{n})$ is Cauchy $\wedge\ (a_{n})$ is a sequence in
$A\Rightarrow(a_{n})$ converges in $A$
Next, it does a move that corresponds closely to the human move of picking an
arbitrary Cauchy sequence: it gets rid of the universal quantifier and places
the premises of the resulting conditional statement above the line, leaving
the conclusion below the line. A small technical point is that we insist that
it combines these two operations into a single move, because we do not allow
“bare” conditionals: that is, statements of the form $P(x)\Rightarrow Q(x)$
that are not universally quantified. (One reason for doing this is that such
statements almost always seem very strange and unnatural to humans. However,
we have also found situations where it helps the program to avoid doing
genuinely bad moves.)
H1. $X$ is a complete space
H2. $A$ is closed in $X$
H3. $(a_{n})$ is Cauchy
H4. $(a_{n})$ is a sequence in $A$
T3. $(a_{n})$ converges in $A$
Next comes an operation that fits under the general heading of _forwards
reasoning_ : from the assumptions that $(a_{n})$ is Cauchy and $X$ is
complete, it follows that $(a_{n})$ converges. Of course, this deduction can
be broken up: one would expand the statement “$X$ is complete” to say that
every Cauchy sequence in $X$ converges, then substitute $(a_{n})$ for the
universally quantified sequence, and finally use modus ponens to deduce that
$(a_{n})$ converges. However, a human mathematician would do all this in one
step, and therefore so does our program. Of course, to do it the program must
have _access_ to the expansion of “$X$ is complete”, which indeed it does: it
has a library of definitions and basic facts that it can make use of whenever
it wants. This models the collection of definitions and facts that a human
mathematician would store in his or her long-term memory.
After the step, the problem state is this (except that there will be tags on
the statements H1 and H3 to indicate that they have been used – to keep the
problem states easy to read we are not displaying the tags in this
discussion).
H1. $X$ is a complete space
H2. $A$ is closed in $X$
H3. $(a_{n})$ is Cauchy
H4. $(a_{n})$ is a sequence in $A$
H5. $(a_{n})$ converges
T3. $(a_{n})$ converges in $A$
At this point, it will be clear to a human that the hypothesis H1 has been
‘used up’. That is, it has played its role and will almost certainly not be
required again. It is not easy to work out precise necessary and sufficient
conditions for when humans judge a statement to be used up in this sense, but
we have identified some conditions that appear to be sufficient. The
hypothesis H1 satisfies those conditions and is therefore deleted.
A human mathematician would also be confident that hypothesis H3 has been used
up. Our program does not delete H3 because it does not know that the Cauchy
assumption will not be expanded later. This we regard as an imperfection of
the program that needs at some point to be corrected. (One possible
justification for deleting H3 is that it is implied by H5. Whether this is the
‘right’ reason remains to be worked out.)
Deleting statements turns out to have no effect on the performance of this
program, so there might seem to be no point in worrying about it. However, we
remain of the view that we should try to model _all_ aspects of human problem
solving. A potential benefit of thinking about the discarding of hypotheses is
that it is a special case of a more general issue that undoubtedly _will_ be
important for performance when the problems get harder: dealing with
irrelevant statements. Often with complicated problems one is presented with
more information than one needs, and it is important to decide which out of
all the statements available are likely to be useful.
H2. $A$ is closed in $X$
H3. $(a_{n})$ is Cauchy
H4. $(a_{n})$ is a sequence in $A$
H5. $(a_{n})$ converges
T3. $(a_{n})$ converges in $A$
The next move applied by the program is to expand hypothesis H5. The reason
the program likes this move is that the expansion of the hypothesis begins
with an existential quantifier, so expanding it gives us a new object to work
with. The program automatically removes the existential quantifier, thereby
implicitly picking a limit for the sequence $(a_{n})$ – again, this
surreptitious removal of quantifiers mirrors the way human mathematicians
think and write. After the expansion and quantifier removal, we arrive at the
following problem state.
H2. $A$ is closed in $X$
H3. $(a_{n})$ is Cauchy
H4. $(a_{n})$ is a sequence in $A$
H6. $(a_{n})\to a$
T3. $(a_{n})$ converges in $A$
Now the program does another kind of forwards reasoning, where a result from
the library is used. The library result states that a closed set contains its
limit points. If we apply this result to the set $A$, the sequence $(a_{n})$
and the limit $a$, then the hypotheses H2, H4 and H5 give us precisely the
premises we need, which allows us to conclude that $a\in A$.
H2. $A$ is closed in $X$
H3. $(a_{n})$ is Cauchy
H4. $(a_{n})$ is a sequence in $A$
H6. $(a_{n})\to a$
H7. $a\in A$
T3. $(a_{n})$ converges in $A$
At this point our deletion rules allow the program to delete H2. We also see
why it is important to be cautious about deleting statements: the hypotheses
H4 and H6 have been used, but they are going to be used again. Fortunately,
our deletion rules do not cause either of them to be deleted.
H3. $(a_{n})$ is Cauchy
H4. $(a_{n})$ is a sequence in $A$
H6. $(a_{n})\to a$
H7. $a\in A$
T3. $(a_{n})$ converges in $A$
The final step for the program is to look at the expansion of the target,
which is
$(a_{n})$ is a sequence in $A\ \wedge\ \exists z\ z\in A\ \wedge\ (a_{n})\to
z$
and observe that if it sets $z$ equal to $a$, then all three resulting
statements occur above the line as hypotheses. Because of this, it declares
the problem solved.
### 2.3 Some terminology
With the possible exception of our rules for deleting statements, which are
not essential to the program, all the move types we use in the program are
fairly standard in the field, or are combinations and minor variations of
standard move types.
Before we say what they are, it will be useful to have some terminology for
the different kinds of statements that can occur. We call a statement _atomic_
if it is of the form $P(x_{1},\dots,x_{k})$, where $P$ is a predicate and
$x_{1},\dots,x_{k}$ are terms. We do not insist that it is actually written in
this form, so for example the statement
$d(x,y)<\epsilon$
is an atomic statement, which we could if we wanted to write in the form
is_less_than$(d(x,y),\epsilon)$.
When a statement is built out of atomic statements using connectives and
quantifiers, we classify it according to the operation that appears at the top
of its parse tree. For example, the statement
$\exists x\ x\in A\ \wedge\ d(x,y)<\epsilon$
is an _existential_ statement, whereas the statement
$x\in A\ \wedge\ d(x,y)<\epsilon$
is _conjunctive_. It is often useful to look more than one level down the
parse tree. For example, we would call the statement
$\forall x\ x\in A\Rightarrow x\in B$
a _universal conditional_ statement. (Recall that we do not allow ‘bare’
conditional statements, so this is a particularly important category.)
Similarly, the existential statement above can be further classified as an
_existential conjunctive_ statement.
Finally, many atomic statements can be expanded into statements that are no
longer atomic. For example, the statement $A\subset B$ expands to the
universal conditional statement above. It is often useful to know what a
statement will become after it is expanded: to specify this we use the prefix
‘pre-’. Thus, the statement $A\subset B$ is pre-universal conditional. We call
an expansion _elementary_ if it does not introduce a quantifier. For example,
the expansion of $A\subset B$ is not elementary, whereas the expansion of
$x\in A\cap B$
as
$x\in A\ \wedge\ x\in B$
is elementary.
### 2.4 How the program works
The following lines are taken directly from the program’s code: they list, in
order of priority, the names of the moves that it can do. In the rest of the
section, we shall explain what those moves are.
--Deletion
deleteDone,
deleteDoneDisjunct,
deleteDangling,
deleteUnmatchable,
--Tidying
peelAndSplitUniversalConditionalTarget,
splitDisjunctiveHypothesis,
splitConjunctiveTarget,
splitDisjunctiveTarget,
peelBareUniversalTarget,
removeTarget,
collapseSubboxTarget,
--Applying
forwardsReasoning,
forwardsLibraryReasoning,
expandPreExistentialHypothesis,
elementaryExpansionOfHypothesis,
backwardsReasoning,
backwardsLibraryReasoning,
elementaryExpansionOfTarget,
expandPreUniversalTarget,
solveBullets,
automaticRewrite,
--Suspension
unlockExistentialUniversalConditionalTarget,
unlockExistentialTarget,
expandPreExistentialTarget,
convertDiamondToBullet,
--EqualitySubstitution
rewriteVariableVariableEquality,
rewriteVariableTermEquality
The program repeatedly applies a move of the first type it can from this list.
Thus, if a move of type deleteDone can be performed, it performs it. If not,
but a move of type deleteDoneDisjunct can be performed, then it performs that.
Otherwise, it tries deleteDangling. And so on.
#### 2.4.1 Deletion moves
The moves can be divided into some broad categories. We begin by discussing
moves that delete statements from problem states.
Deletion is intended to model the human capacity to recognise that a statement
is no longer going to be used and to turn attention away from it. Exactly how
mathematicians do this involves several subtleties, which we shall discuss
elsewhere. Here we give a brief indication of how the program works. One key
point is that we do not normally delete statements unless they have been used
as part of a reasoning step. This is because for the problems our program
tackles, there is a strong presumption that every statement that appears will
at some point be used, so it is risky to delete an unused statement, even if
it appears to be unusable. However, once a statement has been used, this risk
is greatly diminished, so under suitable conditions we are much happier to
delete it. Therefore, the program tags hypotheses as ‘vulnerable’ when they
have been used, and its deletion rules are not applied to hypotheses unless
they have this tag.
#### deleteDone
There are several situations where a move results in a target being replaced
by the word ‘done’ because it has been proved. Once this has happened, the
program immediately deletes it. The aim of the program is to reach a problem
state with no targets.
#### deleteDoneDisjunct
If a target is disjunctive and one of its disjuncts is the word ‘done’, then
the entire target is deleted.
#### deleteDangling
We call a free variable $v$ _dangling_ if it is involved in just one
statement. This is a strong sign that the statement cannot be used. It is not
a conclusive proof: for example, there might be an existential target that can
be solved if we substitute $v$. However, if a statement contains a dangling
variable and is vulnerable, then the program deletes it. This is a pragmatic
decision on our part rather than one that is fully justified theoretically:
there may be problems where deleting statements that have been used and that
contain dangling variables is the wrong thing to do, but that situation does
not appear to arise for the kinds of problems the program is designed to
handle.
#### deleteUnmatchable
Suppose that we have the statements $x\in A$ and $A\subset B$ as hypotheses.
The expansion of $A\subset B$ is $\forall u\ u\in A\implies u\in B$. If we
substitute $x$ for $u$, then the premise of this statement becomes $x\in A$,
which is identical to the hypothesis. We say that $x\in A$ _matches_ the
premise of (the expansion of) $A\subset B$. We call a statement _unmatchable_
if there are no available matches for it.
The program is not allowed to substitute the same variable twice into the same
hypothesis. (This is partly because no human would ever do so, and partly to
avoid getting into a loop.) This can create further circumstances where a
hypothesis is unmatchable. For example, suppose we apply forwards reasoning to
the statements $x\in A$ and $A\subset B$ to deduce that $x\in B$. Then we can
no longer use the match between $x\in A$ and $A\subset B$, so $x\in A$ becomes
unmatchable (assuming that there is no other statement that matches it). Since
it has been used, it is vulnerable, and will therefore be deleted. If no other
statement matches $A\subset B$, then that too will be deleted.
#### 2.4.2 Tidying moves
Tidying moves are moves that do not substantially change the logic of the
problem state, but put it into a more convenient form.
#### peelAndSplitUniversalConditionalTarget
If the target takes the form $\forall x\ P(x)\Rightarrow Q(x)$, then this move
creates a new hypothesis $P(x)$ and replaces the target by $Q(x)$. This
corresponds to the human move of saying (or thinking), ‘Let $x$ be such that
$P(x)$; we need to show that $Q(x)$.’ We can regard it as a composition of two
moves: one to get rid of the quantifier and one to split up the conditional
statement. We use the word ‘peel’ to refer to any move that gets rid of a
quantifier.
If there is more than one target, then this move has to be modified, since we
cannot use $P(x)$ to help us prove a different target. In that situation, we
create what we call a _box_. A box looks like a problem state within a problem
state: it has a line with some statements above it and some statements (or
subboxes – box formation can be iterated) below it. If the peel-and-split move
creates a box, then that box lives below the main line.
That is, if we have a problem state of the form
Hypotheses
$\forall x\ P(x)\implies Q(x)$
$R$
then after the move we will transform it to
Hypotheses
$P(x)$
---
$Q(x)$
$R$
The program then knows that it can use $P(x)$ to prove $Q(x)$ but not to prove
$R$. The hypotheses above the main line can of course be used to prove both
statements.
Boxes are regarded as complex targets (in the above case the target
corresponding to the box would be to demonstrate that $P(x)$ implies $Q(x)$)
and as such can form part of a list of targets. For example, if there are two
boxes, one written above the other, that means that the implications within
both boxes need to be established. Sometimes one needs to establish just one
such implication: in that case one writes the two boxes side by side with a
‘$\vee$’ symbol in between them.
In theory, boxes can be nested, though this rarely happens in practice. The
rules governing which statements can be used to prove what are as follows. For
each statement, there is a minimal box that contains it (counting the entire
problem state as a box). Two boxes are _comparable_ if one contains the other
in an obvious sense, and two statements are _comparable_ if the minimal boxes
containing them are comparable. All deductions that the program makes must
involve comparable statements.
A quick way to interpret a box is to regard the line in the middle as an
implication sign, the lists of statements involved as linked by conjunction,
and the box itself as placing brackets round everything.
#### splitDisjunctiveHypothesis
If there is a hypothesis of the form $P\vee Q$, then the program splits the
problem into two, one with $P$ as a hypothesis and one with $Q$ as a
hypothesis. Technically, it achieves this by forming two boxes, one displayed
above the other to indicate that both must be established.
#### splitConjunctiveTarget
If there is a target of the form $P\wedge Q$, then it is replaced by two
targets $P$ and $Q$.
#### splitDisjunctiveTarget
If there is a target of the form $P\vee Q$, then it is replaced by two boxes,
linked by a $\vee$ symbol. One box has nothing above the line and $P$ below
the line, and the other has nothing above the line and $Q$ below the line.
This move exists for technical reasons: for example, the program sometimes
likes to attach tags to statements, but it has no facility for attaching tags
to parts of statements. Therefore, if we want to use a tag to record
information about one disjunct of a disjunctive target, we need to ‘split’ the
target first.
#### peelBareUniversalTarget
If the target is of the form $\forall x\ P(x)$ and $P$ is not a conditional
statement, then this move replaces the target by $P(x)$.
There is a sense in which bare universal targets should never occur, since
when we quantify over $x$, we do not quantify over every object in the
universe, but rather over some set $X$. So one might argue that the statement
$\forall x\ P(x)$ should really be rendered as the universal conditional
statement $\forall x\ x\in X\implies P(x)$. What we mean by a ‘bare universal’
statement is one where the ‘premises’ are background statements that we do not
want to elevate to ‘substantive’ status. For example, we would regard the
expansion of the statement ‘$G$ is Abelian’ as a bare universal, since the
only condition needed for two elements to commute is the background
information that they are elements of $G$.
#### removeTarget
This is actually a class of move types, but what they have in common is that
under appropriate circumstances they replace a target with the word ‘done’.
The most obvious example is when a target equals a hypothesis (and that
hypothesis is allowed to be used to prove the target). A more complicated
example is when the target is of the form $\exists u\ P(u)\wedge Q(u)$ and
there are hypotheses $P(x)$ and $Q(x)$. The other circumstances are similar.
#### collapseBoxedTarget
If a box B has nothing above its internal line, contains no suspended
variables (see below for a definition) and is not joined to another box with
the $\vee$ symbol, then the statements below its internal line are listed as
targets without the box B.
#### 2.4.3 Applying moves
An applying move is a move where we apply a hypothesis, result or definition.
#### forwardsReasoning
This covers a number of closely related steps that the program can take. The
most basic example is using hypotheses of the form $P(x)$ and $\forall u\
P(u)\Rightarrow Q(u)$ to obtain the hypothesis $Q(x)$. However, there are
several natural variants and generalizations of this. An obvious one is that
if there are hypotheses of the form $P_{1}(x),\dots,P_{k}(x)$ and $\forall u\
P_{1}(u)\wedge\dots\wedge P_{k}(u)\implies Q(u)$, then the program generates
the hypothesis $Q(x)$. Deductions of this kind are often called _forward
chaining_.
A variant that is worth highlighting is illustrated by the following simple
piece of reasoning: if we know that $x\in A$ and that $A\subset B$ then we can
deduce that $x\in B$. Humans will make this deduction in one step rather than
first expanding the statement $A\subset B$ as $\forall u\ u\in A\Rightarrow
u\in B$. Our program does the same. In general, for each type of reasoning
move that involves a universal conditional hypothesis, there is a variant that
does essentially the same thing to a pre-universal conditional hypothesis.
#### forwardsLibraryReasoning
This is reasoning that is ‘mathematical’ rather than ‘purely logical’. We have
seen an example of it already: deducing from the statements ‘$(a_{n})$ is a
sequence in $A$’, ‘$A$ is closed’ and ‘$a_{n}\to a$’ that $a\in A$. The reason
this deduction can be done in one step is that the library contains a general
result that says that whenever a sequence in a closed set tends to a limit,
then the limit belongs to the closed set as well.
Logically speaking, forwards library reasoning is similar to ordinary forwards
reasoning, but there are one or two aspects of it that give it a different
flavour. The main one is that library results contain no free variables: they
are general facts that apply universally. This distinguishes them from
hypotheses, which are more contingent. A second difference is that forwards
library reasoning is normally used to deduce an atomic hypothesis from other
atomic hypotheses. A universal conditional statement is involved, but it is in
the library and is not a hypothesis.
#### expandPreExistentialHypothesis
As its name suggests, this means replacing a pre-existential hypothesis by its
expansion. What the name of the move does not reveal is that this expansion is
followed immediately by a peeling to get rid of the existential quantifier. So
for example the statement ‘$\alpha$ has an inverse’ might be replaced by
$\alpha\beta=\beta\alpha=1$.
This would be modelling a two-step human thought process. The first step is to
note that there exists $\beta$ such that $\alpha\beta=\beta\alpha=1$, and the
second step is quietly to forget about the existential quantifier and to refer
to $\beta$ as though it has been chosen. Human mathematicians will usually
miss out a step that says something like, ‘Let $\beta_{0}$ be such that
$\alpha\beta_{0}=\beta_{0}\alpha=1$’, so the program does as well.
#### elementaryExpansionOfHypothesis
This takes a hypothesis that has an elementary expansion and replaces it by
that expansion. This is sometimes combined with some tidying. For example, if
the hypothesis in question is $x\in A\cap B$, then the elementary expansion is
$x\in A\wedge x\in B$, but this expansion is immediately converted into the
two hypotheses $x\in A$ and $x\in B$ and does not itself appear in any problem
state.
#### backwardsReasoning
Given a target $Q(x)$ and a hypothesis $\forall u\ P(u)\Rightarrow Q(u)$, this
replaces the target by $P(x)$.
More generally, if we have a target $Q(x)$ and a hypothesis $\forall u\
P_{1}(u)\wedge\dots\wedge P_{k}(u)\implies Q(u)$, then it is logically sound
to replace the target $Q(x)$ by the $k$ targets $P_{1}(x),\dots,P_{k}(x)$.
Deductions of this kind are often called _backward chaining_. The program is
allowed to do this more complex backward chaining only under tightly
constrained circumstances: it must be that all but one of the statements
$P_{1}(x),\dots,P_{k}(x)$ is a hypothesis, so that only one new target is
created. This is another pragmatic decision: it is a crude way of deciding
whether applying the hypothesis $\forall u\ P_{1}(u)\wedge\dots\wedge
P_{k}(u)\implies Q(u)$ is likely to be the right thing to do, and the severity
of the test is intended to stop the program making ‘speculative’ deductions
that risk leading to combinatorial explosion.
As with forwards reasoning, there is a simple variant where the role of the
universal conditional hypothesis is played by a pre-universal conditional
hypothesis instead. For example, given a target $x\in B$ and a hypothesis
$A\subset B$ the program could use this variant to replace the target by $x\in
A$.
#### backwardsLibraryReasoning
This is backwards reasoning that makes use of a general result in the library.
However, it is slightly subtler than forwards library reasoning, because it
always uses hypotheses as well as a target. The precise rule is that if there
are hypotheses $P_{1}(x),\dots,P_{k-1}(x)$, a library result $\forall u\
P_{1}(u)\wedge\dots\wedge P_{k}(u)=Q(u)$ and a target $Q(x)$, then the target
can be replaced by $P_{k}(x)$. (The premises of the library result do not have
to be stated in the order $P_{1},\dots,P_{k}$.)
An example of this kind of reasoning would be to say, “It is sufficient to
prove that $B$ is open,” if one wished to prove that $A\cap B$ was open and
knew that $A$ was open. This would be making use of the result that an
intersection of two open sets is open.
#### elementaryExpansionOfTarget
This replaces a target by an elementary expansion of that target, if it has
one.
#### expandPreUniversalTarget
This replaces a pre-universal target by its expansion. This move will be
followed by one of the tidying moves peelAndSplitUniversalConditionalTarget or
peelBareUniversalTarget. It is usually the first move that the program makes
when faced with a naturally stated problem.
#### solveBullets
As we are just about to discuss in more detail, we sometimes mark a variable
$w$ with a diamond or a bullet. This indicates that the variable needs at some
stage to be chosen in such a way that the problem can be solved. If the
variable only ever appears in targets, then one simple way in which this can
often be done is to identify another variable $x$ with the property that if we
substitute $x$ for $w$, then every target that involves $w$ is equal to a
hypothesis. In that situation, all those targets are replaced by ‘done’. This
move is what we call ‘solveBullets’.
#### automaticRewrite
There are a few rewriting rules stored in the library. Two examples are that
the statement $x\in f^{-1}(A)$ is rewritten as the statement $f(x)\in A$ and
the term $g\circ f(x)$ is rewritten as the term $g(f(x))$. (Of course, these
are general rewriting rules and work whatever the variables happen to be
called.)
The rewriting of statements takes place only when those statements have been
isolated as hypotheses or targets, so for example the program would not
rewrite the statement $y\in f^{-1}(U)$ when it occurs inside the larger
statement $\forall y\ d(x,y)<\delta\Rightarrow y\in f^{-1}(U)$. However, terms
can be rewritten as soon as they appear.
#### 2.4.4 Suspending moves
We now come to the class of moves just alluded to: moves that help us deal
with existential targets when it is not immediately clear what to substitute
for the existentially quantified variable. A standard technique for this,
which is essentially the technique we use, is to form _metavariables_. The
rough idea of a metavariable is that one reasons with it as though it had been
chosen, deferring the actual choice until later when it becomes clearer what
choice will make the argument work. Mathematicians often use this trick: a
classic example is the ‘$3\epsilon$-argument’ used to prove that a uniform
limit of continuous functions is continuous.
When the program ‘pretends that it has chosen’ a variable, it marks that
variable with a bullet, and we say that it has been _suspended_. Thus,
suspension is the process of converting a variable into a metavariable.
However, we found it convenient to introduce two ‘levels of suspension’, to
model two styles of reasoning that are logically similar but psychologically
quite different.
#### unlockExistentialUniversalConditionalTarget
To illustrate this, suppose we have a target such as $\exists\delta\ \forall
y\ d(x,y)<\delta\Rightarrow f(y)\in B$, and also a hypothesis $\forall u\ u\in
A\implies f(u)\in B$. Then it is easy to see that we can reduce the target to
$\exists\delta\ \forall y\ d(x,y)<\delta\implies y\in A$. However, this move
is not open to the program because it is not allowed to ‘reason inside
quantifiers’. This is a matter of convenience: such moves are logically valid,
but it is tedious to specify appropriate variants of several of the reasoning
moves listed above. Instead, we apply a procedure that we call _unlocking_ ,
which effectively moves aside the existential quantifier and allows the
program to reason as normal with the statements inside it.
More precisely, what the program does to ‘unlock’ the statement is create a
box. In the example above, it would have no statements above the line, and
below the line it would have the statement $\forall y\
d(x^{\blacklozenge},y)<\delta\Rightarrow f(y)\in B$. The diamond on the
variable $x$ indicates that $x$ needs to be chosen.
It is important for the program not to interchange quantifiers accidentally.
For this reason, we tag the box just created with the variable
$x^{\blacklozenge}$, to indicate that the existential quantification over $x$
is within that box.
After unlocking the statement, the program peels and splits the resulting
universal conditional target (so a more accurate name for the move type would
be unlockPeelAndSplitExistentialUniversalConditionalTarget). After that, we
have a box that looks like this.
$d(x^{\blacklozenge},y)<\delta$
---
$f(y)\in B$
Once we have done this, the statement $f(y)\in B$ has become a target and the
program is free to apply backwards reasoning to it.
#### unlockExistentialTarget
This move replaces a target of the form $\exists x\ P(x)$ with a box that has
nothing above the line and the statement $P(x^{\blacklozenge})$ below the
line. The box is labelled with the variable $x^{\blacklozenge}$.
This move will never be applied to an existential universal conditional
target, since that will have been dealt with by
unlockExistentialUniversalConditionalTarget. The main reason we have two
separate moves here is that we prefer to bundle the unlocking together with
peeling and splitting when that is possible.
To see what unlockExistentialTarget allows the program to do, suppose that we
have a target of the form $\exists x\ Q(x)\ \wedge\ R(x)$ and also a
hypothesis of the form $\forall u\ P(u)\Rightarrow Q(u)$. In this situation we
would like to be able to do backwards reasoning inside the existential
quantifier to reduce the target to $\exists x\ P(x)\ \wedge\ R(x)$. However,
the program does not have a move for this. Instead, it unlocks the existential
target, so that it has a box with the statement $Q(x^{\blacklozenge})\ \wedge\
R(x^{\blacklozenge})$ below the line. The tidying move splitConjunctiveTarget
can now turn this new target into two targets, and once it has done that, the
applying move backwardsReasoning can be used to replace the target
$Q(x^{\blacklozenge})$ by $P(x^{\blacklozenge})$.
As another example of the use of unlocking, suppose that we wished to prove
that $A\cap B$ is non-empty and had the hypotheses $x\in A$ and $x\in B$. The
program cannot see that $x$ is a witness to the non-emptiness of $A\cap B$
without doing some processing. An obvious first step is to expand the target
into the statement $\exists u\ u\in A\cap B$. However, the program is not then
allowed to do an elementary expansion inside the quantifier. Instead, it
unlocks $u$ so that there is a new target $u^{\blacklozenge}\in A\cap B$. This
can now be expanded and split into the two targets $u^{\blacklozenge}\in A$
and $u^{\blacklozenge}\in B$, which solveBullets can then match with the
hypotheses.
This may seem a little circuitous, but it actually models quite closely how
humans think. A human might say, ‘I want to show that $A\cap B$ is non-empty,
so I need to find some $u$ that belongs to $A\cap B$. In other words, I need
$u$ to be in $A$ and in $B$. Aha, I can take $x$.’ The program’s unlocking
models the silent disappearance of the existential quantifier before the
second sentence of the above.
#### expandPreExistentialTarget
This does exactly what it says: it replaces a pre-existential target by its
expansion.
#### convertDiamondToBullet
There are certain moves that the program will not do with a ‘diamonded’
variable. In particular, it will not do any reasoning with a hypothesis that
involves such a variable: for that it needs a deeper level of suspension,
roughly speaking corresponding to the human move of ‘pretending that a
variable has been chosen’ and then reasoning with it. Logically this is not an
important difference, but it is a useful one for us because it reflects a
difference in the way human mathematicians think and write. This helps the
program to produce more convincing write-ups.
We do not need separate move types for reasoning that involves hypotheses with
bulleted variables: we just allow the reasoning moves above to take such
variables.
An important technicality is that if we postpone the choice of a variable, we
must keep track of what other variables it is allowed to depend on. However,
what we actually do is note which variables it is _not_ allowed to depend on.
This is for two reasons. First, it seems to reflect more accurately how human
mathematicians think about such variables, and secondly, it is more
economical: there are typically many fewer variables on which a bulleted
variable is not allowed to depend than variables on which it is allowed to
depend.
#### 2.4.5 Equality substitution
If we are told that two objects are equal, then we can eliminate all mention
of one object in favour of the other. The precise rules governing when and how
mathematicians tend to avail themselves of this opportunity are not obvious.
The rules below are best regarded as a temporary solution: they do not always
result in realistically human choices, and we intend to replace them by more
satisfactory rules in the near future.
#### rewriteVariableVariableEquality
If there is a hypothesis of the form $x=y$, then this move replaces all
occurrences of $y$ by $x$ and eliminates the hypothesis.
#### rewriteVariableTermEquality
If there is a hypothesis of the form $v=t$ or $t=v$, where $v$ is a variable
and $t$ is a term, then this move replaces all occurrences of $t$ by $v$.
### 2.5 Justification for the order of priority
From examining how humans solve simple problems in the theory of metric
spaces, it is not too hard to arrive at the above list of move types. But it
is less clear what principles should govern the architecture – that is, the
way that the program decides which move type to do in any given situation. One
obvious method of choosing an architecture is to work through a large number
of problems and try to observe what seems to be the natural approach. After a
while, one can make a guess at how the program should work, and if the guess
results in strange behaviour for some problems then one can refine it, hoping
that the process of refinement will stabilize quickly. An alternative method
is to try to devise a theory that explains which move is the best than in each
situation: ideally, that will turn out to be the move that humans are
naturally drawn to. The second method is harder, but if it works, then the
advantage is that the rules are likely to be more robust: without a proper
theoretical backing, one cannot be as confident that they will not lead to
inappropriate behaviour when the program is presented with an unfamiliar
problem.
We have used a mixture of the two methods. We have good reasons for some of
the choices we have made, but other choices are justified by the fact that
they seem to work (in the sense of leading to human-like behaviour). A broad
overarching principle that gives a theoretical backing to many of our choices
is this: _the program prefers safe moves to dangerous moves_. The picture we
have here is one where at any stage there is a choice of moves that can be
made, and we have to make an assessment of how likely any given choice is to
form part of the argument one is looking for. The greater this likelihood, the
safer the move.
Because it seems hard to attach probabilities to statements in this way, we
have not tried to do so. However, since something like the above picture seems
to be what humans do, we bear the picture in mind when planning the program.
In particular, if a move is obviously safe, we will assign it a high priority.
A good example of a safe move is a tidying move. If, for example, we have a
conjunctive hypothesis, then there is nothing to lose by splitting it up into
its conjuncts, so that move we do automatically without any hesitation at all.
By contrast, expanding a definition is substantially less safe: sometimes it
is possible to reason in a high-level way without expanding, and since we do
not allow ‘de-expansion’ in this program (and in general allowing it would be
highly problematic because of the danger of an infinite loop), expanding a
definition is closing off the option of such high-level arguments. As an
example, in the problem we discussed earlier, the program does not expand the
statements ‘$(a_{n})$ is Cauchy’, ‘$(a_{n})$ is a sequence in $A$’ or
‘$(a_{n})\to a$’. That allows it to do some high-level forwards library
reasoning that would no longer be possible if any one of those three
statements was expanded.
Thus, expansion has a fairly low priority. Having said that, some expansions,
such as elementary expansions or expansions of pre-existential hypotheses, are
considerably safer, so those ones have higher priority.
Somewhere in between are the other reasoning moves. Here it becomes more
complicated to apply the general principle (even as an informal guiding
principle), because the safety of a move type depends heavily on context. In
particular, forwards reasoning is in general fairly unsafe – if you have a lot
of information and do not know which statements are relevant, then the
probability that any given deduction will form part of the eventual proof may
be quite small – but is much safer when it comes to routine problems, which
tend not to suffer from the problem of irrelevant information.
The psychology literature suggests that _when it is safe_ , humans tend to
prefer forwards reasoning to backwards reasoning swelleretal1983 ;
owensweller1985 , though this appears to be a question more of style than of
problem-solving efficacy: we seem to prefer not to keep track of a moving
target if we do not have to. Since forwards reasoning tends to be safe for the
highly routine problems our program tackles, we have given all forwards
reasoning a higher priority than all backwards reasoning. This also has the
beneficial effect of making the program reluctant to switch direction – too
much switching from forwards to backwards or vice versa would again be bad
mathematical style.
This aspect of our program is, however, unstable, in that we know that in
order to develop the program we will have to change it. In fact, we even have
an example of a rather routine problem where our program performs badly for
this very reason. If it is asked to show that the intersection of two
subgroups $H$ and $K$ of a group $G$ is itself a subgroup of $G$, then when it
is trying to prove that $H\cap K$ is closed under multiplication, it expands,
peels and splits the target, obtaining two elements $x$ and $y$ of $H\cap K$
and a target of showing that $xy\in H\cap K$. At that point, there is nothing
to stop the program making ‘silly’ deductions such as that $x^{-1}\in H$.
There are various easy ways of dealing with this problem and we shall
implement these in future versions: for example, we could add restrictions on
creating new terms (a human would not think of deducing that $x^{-1}\in H$
because ‘we are not interested in $x^{-1}$’) or we could alter the priorities
so that when there are a number of possible forwards moves, so that it is no
longer clear that they are all relevant, the program switches to backwards
reasoning.
One other feature of the ordering of reasoning moves is that we prefer pure
reasoning moves to library reasoning moves. That is because in general a
hypothesis is more likely to be relevant than a library statement, though if
enough of the premises of a library statement are present as hypotheses, that
is a fairly strong argument for its relevance.
At the bottom of the list of priorities is suspension. That is because humans
tend to regard it as a last resort. When mathematicians need to prove
statements of the form $\exists x\ P(x)$, then by and large they prefer to
simplify the problem until a suitable candidate $x_{0}$ for $x$ becomes
obvious and it remains to carry out the relatively easy task of verifying that
$P(x_{0})$. Only when this straightforward approach fails do we take the more
drastic step of pretending that $x$ has been chosen.
We will not say much more here about how we chose the priority order, but we
have two brief further points. First, although our reasons are not 100%
precise, we found that in practice they were adequate, in the sense that they
suggested an order before we started, and we found that we did not have to
modify the order when we tried further problems (though, as commented above,
there are certain aspects of the architecture that will need to be changed in
future versions). Secondly, when it comes to the finer detail of the ordering,
there may not be that much to choose between different move types. However,
conflicts rarely arise between different move types that are not distinguished
by any of the above criteria, so in practice these finer details have little
if any effect on what the program actually does.
### 2.6 Example: An intersection of two open sets is open
Now that we have discussed how the program works, let us look at another
example, which illustrates most of the move types and shows how the order or
priority works in practice. The problem to be solved is the following.
###### Problem 3
Let $A$ and $B$ be open subsets of a metric space $X$. Prove that $A\cap B$ is
open.
The initial problem state is as follows.
$A$ is open
$B$ is open
$A\cap B$ is open
No reasoning moves are possible, so we end up having to expand. The highest
priority move we can do is expandPreUniversalTarget, which, after the tidying
move peelAndSplitUniversalConditionalTarget, has the following effect.
$A$ is open
$B$ is open
$x\in A\cap B$
$\exists\delta\ \forall y\ d(x,y)<\delta\Rightarrow y\in A\cap B$
It may look a little strange that we do not specify that $\delta>0$. The
reason for that is that we think of positivity as a ‘background condition’
rather than as a ‘substantive statement’. One reason that matters is connected
with safety: if a premise of a universal conditional hypothesis is satisfied
by some variable, then that is good evidence that the forwards reasoning one
can do as a result is relevant to the problem. Or rather, it is good evidence
if the premise is a ‘substantive statement’, but not if it is merely a
background condition such as positivity. So the distinction is a useful one.
We think of background conditions as similar to, but not the same as, type
declarations. The program does, however, know that $\delta$ is a positive real
number: that fact is stored as background information when the statement
‘$A\cap B$ is open’ is expanded.
At this point, the reluctance of the program to suspend $\delta$ means that it
does as much forwards reasoning as it possibly can. It begins with
elementaryExpansionOfHypothesis, applied to the third hypothesis.
$A$ is open
$B$ is open
$x\in A$
$x\in B$
$\exists\delta\ \forall y\ d(x,y)<\delta\Rightarrow y\in A\cap B$
This allows it apply two forwardsReasoning twice. After the first application,
the problem state is as follows.
$A$ is open
$B$ is open
$x\in A$
$x\in B$
$\forall u\ d(x,u)<\eta[x]\Rightarrow u\in A$
$\exists\delta\ \forall y\ d(x,y)<\delta\Rightarrow y\in A\cap B$
Note that the last hypothesis is in a sense generated by a combination of
submoves: the first is forwardsReasoning (using the hypotheses $x\in A$ and
‘$A$ is open’) and the second is a peeling (to get rid of $\exists\eta$ at the
beginning of the statement). However, the latter is so automatic that it is
not listed as one of our tidying moves: instead, it is considered as part of
any other move that potentially generates an existential hypothesis.
It is important to keep track of the fact that $\eta$ depends on $x$, which is
what is signified by $\eta[x]$.
After this, deleteUnmatchable causes the program to delete the statements
$x\in A$ and ‘$A$ is open’. This is because both statements have been used, so
they are vulnerable, and because it is no longer permissible to substitute $x$
into ‘$A$ is open’. The resulting problem state is as follows.
$B$ is open
$x\in B$
$\forall u\ d(x,u)<\eta[x]\Rightarrow u\in A$
$\exists\delta\ \forall y\ d(x,y)<\delta\Rightarrow y\in A\cap B$
It then runs through a similar process for $B$ (it does not yet have the
capacity to recognise that the problem is symmetric in $A$ and $B$ and say,
‘Similarly …’). After that process, it arrives at the following.
$\forall u\ d(x,u)<\eta[x]\Rightarrow u\in A$
$\forall v\ d(x,v)<\theta[x]\Rightarrow v\in B$
$\exists\delta\ \forall y\ d(x,y)<\delta\Rightarrow y\in A\cap B$
It has now reached the point where it must suspend $\delta$. In the first
instance, it uses the move unlockExistentialUniversalConditionalTarget, which
includes a peeling and splitting. The result is as follows.
$\forall u\ d(x,u)<\eta[x]\Rightarrow u\in A$
$\forall v\ d(x,v)<\theta[x]\Rightarrow v\in B$
$d(x,y)<\delta^{\blacklozenge}[\overline{y}]$
---
$y\in A\cap B$
The notation $\delta^{\bullet}[\overline{y}]$ signifies that $\delta$ is not
allowed to depend on $y$.
The highest priority move the program can do now is
elementaryExpansionOfTarget, so it does that, and automatically splits the
resulting conjunctive statement (rather than using the move
splitConjunctiveTarget).
$\forall u\ d(x,u)<\eta[x]\Rightarrow u\in A$
$\forall v\ d(x,v)<\theta[x]\Rightarrow v\in B$
$d(x,y)<\delta^{\blacklozenge}[\overline{y}]$
---
$y\in A$
$y\in B$
This allows it to apply backwardsReasoning twice. After the two deductions it
reaches the following state. (It does them separately, so we are jumping a
step here.)
$\forall u\ d(x,u)<\eta[x]\Rightarrow u\in A$
$\forall v\ d(x,v)<\theta[x]\Rightarrow v\in B$
$d(x,y)<\delta^{\blacklozenge}[\overline{y}]$
---
$d(x,y)<\eta[x]$
$d(x,y)<\theta[x]$
It then uses deleteUnmatchable to delete the two hypotheses it has just used.
$d(x,y)<\delta^{\blacklozenge}[\overline{y}]$
---
$d(x,y)<\eta[x]$
$d(x,y)<\theta[x]$
At this point, there is not much that the program can do, because it is not
allowed to reason with the diamonded variable $\delta^{\blacklozenge}$. So the
highest-priority move it can do is convertDiamondToBullet. Also, since there
are no hypotheses above the main line, it may as well remove that line and the
box.
$d(x,y)<\delta^{\bullet}[\overline{y}]$
$d(x,y)<\eta[x]$
$d(x,y)<\theta[x]$
Now it applies backwardsLibraryReasoning. The result in the library is that if
$a<b$ and $b\leq c$, then $a<c$. Applying that with the hypothesis and the
first target results in the following problem state.
$d(x,y)<\delta^{\bullet}[\overline{y}]$
$\delta^{\bullet}[\overline{y}]\leq\eta[x]$
$d(x,y)<\theta[x]$
The deletion rules do _not_ allow the program to delete the hypothesis we have
just used (and this is a good example of a situation where deletion would be a
very bad idea). However, it is aware that it cannot use the hypothesis with
the new target. (We shall not describe here the precise mechanism by which it
gains this awareness.) The program then uses backwardsLibraryReasoning again
and this time it does delete the hypothesis, on the grounds that the variable
$x$ that appears in the hypothesis is dangling. After that, it has reached the
following state.
$\delta^{\bullet}[\overline{y}]\leq\eta[x]$
$\delta^{\bullet}[\overline{y}]\leq\theta[x]$
This is a ‘standard’ existence problem, so the solution is in a library of
standard solutions and the program declares the problem solved. It is here
that the background information that $\delta$, $\eta$ and $\theta$ are
positive is used, since the library result is that the minimum of two positive
real numbers $a$ and $b$ is a positive real number that is less than or equal
to both $a$ and $b$.
That concludes the discussion of this example. For more examples of the
program’s output, as well as the precise output for this problem (which we
have very slightly simplified in one or two places to make the presentation
clearer) see http://people.ds.cam.ac.uk/mg262/robotone.pdf .
### 2.7 Writing up
In general, natural language generation is a complex process. It involves
multiple levels of planning, which draw on both domain knowledge and models of
the intended audience, and also a phase when the actual text is generated,
which draws on syntactic, morphological and lexical information. An overview
of the process may be found in reiter2000building . Because of this
complexity, building a fully fledged natural language generation system is a
major task. Furthermore, since mathematics contains not just English words but
also a large array of distinctive symbols used in distinctive ways, it is not
at all straightforward to use off-the-shelf systems.
Fortunately, mathematical language has properties that make the task
considerably simpler than it is for the English language in general. Foremost
among these is the fact that mathematical proofs almost always have a
particularly simple rhetorical structure. To some degree this is because the
domain of discourse includes only timeless facts, which itself rules out a
large proportion of the rhetorical relations found in general text. But the
main reason is that there is a strong convention that further constrains the
rhetorical structure of proofs. A proof proceeds by the presentation of a
sequence of assertions, each of which follows from the premises of the theorem
being proved or from previous assertions. This structure is not accidental; it
is a direct reflection of the fact that mathematicians process proofs by
reading and verifying one sentence at a time, and would not expect the
justification of a fact presented in one sentence to be deferred to a later
sentence. (We are talking here about proofs of the level of simplicity of the
proofs discussed in this paper. For more complicated arguments, facts may
sometimes be used before they have been proved, but in good mathematical
writing this will be carefully flagged up to make it as easy as possible for
the reader to check that the resulting argument is complete and not circular.)
This convention gives us an easy way to produce write-ups of our proofs. An
obvious strategy is to allow each move to generate some number of sentences
(possibly zero), and to concatenate output from the different moves to produce
the final text. Note that this strategy is viable only because we are
absolutely rigorous about requiring our moves to reflect steps in human
reasoning; in effect, the strategy is mimicking a human who is carefully
writing down down a proof while coming up with it, which is quite
straightforward for an experienced mathematician. (Again, this becomes less
true if the proofs are more difficult.) As we shall see below, this simple
strategy produces surprisingly good results, though with a weakness that needs
to be dealt with by a postprocessing phase, which turns out to be
straightforward.
Because we have a fixed list of move types, implementing the strategy only
requires us to specify which sentences (if any) are produced for the moves of
each type. A very simple way to do this is to use _template generation_ : each
move type is associated with a _template_ , or ‘piece of text with holes’, and
the holes are filled in with concrete information about the facts and objects
used in the actual move performed. So, for example, forwards reasoning may be
associated with a very simple template ‘since $<\textit{facts}>$,
$<\textit{deduced fact}>$’. Instantiating this template would produce text
like
> Since $A$ is open and $x\in A$, there exists $\eta>0$ such that $u\in A$
> whenever $\textit{d}(x,u)<\eta$.
Note that individual facts are expressed in idiomatic ways, rather being
displayed in a way that directly reflects the underlying predicate calculus;
thus we have ‘$A$ is open’ and ‘$\eta>0$’ rather than ‘$\textit{open}(A)$’ and
‘$\textit{greater\\_than}(\eta,0)$’. The same is true of objects: we display
‘$f\circ g$’ rather than compose(f,g), and so on. Similarly quantification is
expressed idiomatically using words like ‘whenever’, where possible, rather
than using more stilted phrases like ‘for all’, which would more directly
reflect the underlying predicate calculus.
An example of the text produced by this method is as follows:
> Let $x$ be an element of $A\cap B$. Since $x\in A\cap B$, $x\in A$ and $x\in
> B$. Since $A$ is open and $x\in A$, there exists $\eta>0$ such that $u\in A$
> whenever $d(x,u)<\eta$. Since $B$ is open and $x\in B$, there exists
> $\theta>0$ such that $v\in B$ whenever $d(x,v)<\theta$. We would like to
> find $\delta>0$ s. t. $y\in A\cap B$ whenever $d(x,y)<\delta$. But $y\in
> A\cap B$ if and only if $y\in A$ and $y\in B$. We know that $y\in A$
> whenever $d(x,y)<\eta$. We know that $y\in B$ whenever $d(x,y)<\theta$.
> Assume now that $d(x,y)<\delta$. Since $d(x,y)<\delta$, $d(x,y)<\eta$ if
> $\delta\leq\eta$. Since $d(x,y)<\delta$, $d(x,y)<\theta$ if
> $\delta\leq\theta$. We may therefore take $\delta=\min\\{\eta,\theta\\}$. We
> are done.
The main problem with this text is that it suffers a lack of coherence, in the
sense defined in knott1996thesis : the sentences are individually acceptable,
but they do not combine to form an idiomatic discourse. The principal reason
for this is that the text repeats information unnecessarily. For example, in
> Since $x\in A\cap B$, $x\in A$ and $x\in B$. Since $A$ is open and $x\in A$,
> there exists $\eta>0$ such that $u\in A$ whenever $\textit{d}(x,u)<\eta$.
the repetition of the underlined phrase is awkward. Because it is introduced
by the sentence immediately preceding the ‘since’ clause, it is awkward to
have it spelt out explicitly within that clause. Similarly, consider:
> Since $\textit{d}(x,y)<\delta$, $\textit{d}(x,y)<\eta$ if $\delta\leq\eta$.
> Since $\textit{d}(x,y)<\delta$, $\textit{d}(x,y)<\theta$ if
> $\delta\leq\theta$.
Here, having two identical ‘since’ clauses in consecutive sentences is again
awkward: the repetition of material is unwieldy and unidiomatic.
We are of the opinion that knott1996thesis correctly diagnoses the underlying
problem here. Spelling out rhetorical relations, or aspects of rhetorical
relations, that can easily be inferred from the context violates Grice’s maxim
of quantity grice1975logic . Often the solution is to substitute an
appropriate and less explicit _cue phrase_. For example, ‘since $A$ is open
and $x\in A$, …’ is better replaced by ‘therefore, since $A$ is open, …’. The
cue phrase ‘therefore’ (which assume the relevant reason has just been given)
is less explicit than the cue phrase ‘since’ (which subordinates an explicitly
stated reason), so it avoids spelling out information that is clear from the
context. In other cases repetition can be avoided by combining sentences; for
example, the second example may be changed into
> Since $\textit{d}(x,y)<\delta$, $\textit{d}(x,y)<\eta$ if $\delta\leq\eta$
> and $\textit{d}(x,y)<\theta$ if $\delta\leq\theta$.
The initial ‘sentence by sentence’ process described above is followed by a
series of transformations that manipulate pairs of consecutive sentences in
order to resolve the issues just mentioned. (Needless to say, the
transformations operate on a structural level rather than on the literal
text.) Applying this series of transformations to the example text above
yields:
> Let $x$ be an element of $A\cap B$. Then $x\in A$ and $x\in B$. Therefore,
> since $A$ is open, there exists $\eta>0$ such that $u\in A$ whenever
> $d(x,u)<\eta$ and since $B$ is open, there exists $\theta>0$ such that $v\in
> B$ whenever $d(x,v)<\theta$. We would like to find $\delta>0$ s. t. $y\in
> A\cap B$ whenever $d(x,y)<\delta$. But $y\in A\cap B$ if and only if $y\in
> A$ and $y\in B$. We know that $y\in A$ whenever $d(x,y)<\eta$ and that $y\in
> B$ whenever $d(x,y)<\theta$. Assume now that $d(x,y)<\delta$. Then
> $d(x,y)<\eta$ if $\delta\leq\eta$ and $d(x,y)<\theta$ if $\delta\leq\theta$.
> We may therefore take $\delta=\min\\{\eta,\theta\\}$ and we are done.
One particular point worth emphasising is that the write-up process is
deterministic: it will always produce the same output text for any given
proof. This is for two reasons. First, if any non-determinism had been present
we would have had to evaluate many outputs for any given proof, which would
have made iterative improvement and fine-tuning of the write-ups considerably
slower. Secondly, and more importantly, if the process were nondeterministic,
our claim that the program produced human-like output would be suspect, in
that we would have been able to run the program several times and ‘cherry
pick’ output. Unfortunately, this determinism has an undesirable (but fully
anticipated) side-effect. When one compare several proofs produced by the
program, the write-ups are much more similar than those a human would produce.
For example, most proofs produced by the program end with the phrase ‘we are
done’. In the long run, we will undoubtedly need to introduce nondeterministic
stylistic variation, allowing the program to vary the text generated for a
particular step in just the way human would, despite the difficulties that
will cause.
Finally, it is worth noting that during the evaluation process described in
the next section, we collated a wealth of data on how humans write up proofs.
We anticipate using this data in combination with carefully chosen natural
language processing techniques to create substantially improved versions of
the write-up procedure.
## 3 Testing the write-ups
Once the program had generated the write-ups for several problems, we wanted
to test whether they could pass for write-ups written by a human
mathematician. In this section we describe an informal experiment that we
carried out for this purpose.
We began by asking two mathematicians, one an undergraduate and one a PhD
student, to write out proofs for five problems for which our program had
generated proof write-ups. We did not tell either of them why we were making
this unusual request, and we did not ask them to make their write-ups as good
as possible. One of the problems was to show that the inverse image of an open
set under a continuous function is open, and one of our volunteers decided to
prove the converse, so that he could use the topological definition of
continuity to prove another of the assertions – that a composition of
continuous functions is continuous. We had to ask him to rewrite the latter
and give the epsilon-delta proof, since we wanted the differences between the
write-ups to be a matter of style rather than substance.
We had another problem of this kind, which was that both our volunteers made
frequent use of open balls. For example, their expansion of ‘$A\cap B$ is
open’ was ‘for every $x\in A\cap B$ there exists $\delta>0$ such that
$B_{\delta}(x)\subset A\cap B$.’ This made some of their arguments neater than
the ones produced by our program. We contemplated getting the program to redo
the problems using open-balls expansions, but in the end decided that it would
be ‘cheating’ to make changes to its output in response to the human write-ups
we had solicited, so we left things as they were.
The program’s write-ups were not designed to be indistinguishable from human
write-ups: we merely wanted them to be acceptable as human write-ups.
Therefore, we left in certain features, such as ending every proof with the
words, ‘we are done’, that we could with a little trouble have changed. (See
the brief discussion of non-determinism at the end of the previous section.)
For this reason, we did not want to ask people to guess which write-ups were
the by the program. Instead, we presented all fifteen write-ups – two by
humans and one by the program for each of the five problems – on the second
author’s blog, and asked readers of the blog to comment on them in any way
they liked. We also asked them to award points for clarity and style. The
orders of the write-ups were chosen randomly and independently. (The precise
mechanism was to decide on a one-to-one correspondence between the set
$\\{1,2,3,4,5,6\\}$ to the set of permutations of the set $\\{1,2,3\\}$, then
to find a website that produced random dice rolls.) So that answers would be
as independent as possible, all comments and ratings were sent to the blog’s
moderation queue and published only after the experiment was finished and
comments on the blog post were closed.
The post can be found at http://gowers.wordpress.com/2013/03/25/an-experiment-
concerning-mathematical-writing/ , together with all the comments and ratings,
but the real point of the experiment was to see whether anybody noticed that
not all the write-ups were by humans. Nobody expressed the slightest suspicion
of this kind.
Having said that, we should also remark that many commenters were highly
critical of the program’s output. Three criticisms in particular stand out.
First, as we expected, the fact that the program did not use open balls was
unpopular: many people commented that this made the write-ups unwieldy.
Secondly, several of the human write-ups stated the new target when the
initial one had been expanded, peeled and split. Several readers commented
that they found this helpful, and criticized our program for not doing it. And
thirdly, commenters did not like the way the program spelt out in detail how
it thought of the right variable to substitute into existential targets (such
as choosing $\min\\{\eta,\theta\\}$ for $\delta$ in the intersection-of-open-
sets problem.
It would be easy to modify the program so that none of these criticisms apply,
so they do not point to fundamentally non-human aspects of how it thinks. To
change the first, we would just have to use open-balls expansions of
definitions such as ‘$A$ is open’ and ‘$f$ is continuous’. To change the
second, we could alter the rule for what the write-up does when we expand,
peel and split, so that it states the new target (preceded by a phrase such as
‘We need to show that’). The third criticism would be harder to deal with, but
when we switch to having two styles of write up – a ‘proof write-up’ and a
‘proof-discovery account’ – then for the first style we will let the program
work out the values of bulleted variables, then simply declare those values
when the variable is first mentioned after being suspended. This will
correspond to the human practice of writing something like ‘Let
$\delta=\min\\{\eta,\theta\\}$’ or ‘Consider the sequence $(b_{n})$ defined by
$b_{n}=a_{n}/(1+a_{n})$,’ which ‘magically’ does exactly what it needs to do
later in the proof.
Although our program’s output came in for quite a bit of criticism, so did the
write-ups by the undergraduate and PhD student – it seems that the readers
were harsh judges. However, for most of the problems, the human write-ups were
found preferable to the program’s.
After the success (as we considered it) of this experiment, we dared to try a
direct test. We published a new post, this time explaining that one proof was
by a program, one by an undergraduate and one by a PhD student, and inviting
readers to vote on which one they thought was by the program. For each
problem, the write-ups were numbered (a), (b) and (c). There were seven
options for the voting: one could choose between (a), (b) and (c), but also
choose between ‘The computer-generated output is definitely $*$’ and ‘I think
the computer-generated output is $*$ but am not certain’; the seventh option
was ‘I have no idea which write-up was computer generated.’ Again there was
the opportunity to comment, for those who wanted to explain the reasons for
their choices.
We did not reveal the results of the voting so far, or anybody’s comments,
until the experiment was ended and the voting was closed. However, there was a
different kind of dependence between answers, which was that people had the
opportunity to look for clues that two different write-ups were from the same
source. Given that we had not tried to remove stylistic ‘tics’ from our
program’s write-ups, this put the program at a significant disadvantage. It
was clear from the comments that many people had noticed that for each problem
exactly one write-up ended with the words ‘we are done’.
Despite this, the program did reasonably well at fooling people that it was
human. The typical pattern was that roughly half the voters would correctly
guess which output was by the program, with slightly under half of that half
saying that the output was definitely by the program. The undergraduate would
always ‘come second’, and there would always be a fair number of people who
said that they had no idea which output was written by the computer. There
were surprisingly many votes for ‘The computer-generated output is definitely
$*$,’ when $*$ was the wrong answer. The total number of votes was always at
least 300, and for the first problem listed (the intersection of open sets is
open) it was over 1000. One slight complication was that after a day or two
the post was listed on the front page of Hacker News. The result was that the
number of votes doubled in a couple of hours, and it may be that the profile
of a typical voter changed. Fortunately, we had noted down the voting numbers
just before this happened, so we presented those results as well as the final
numbers. In the end, however, the proportions did not change very much. The
detailed numbers can be found here:
http://gowers.wordpress.com/2013/04/14/answers-results-of-polls-and-a-brief-
description-of-the-program/ .
One thing this experiment could not tell us, except to a limited extent
through the comments, was whether the program was good at fooling
_mathematicians_ that it was human. It could be that the more mathematically
experienced readers found the program’s output easy to distinguish, while the
votes for the human write-ups came from people who were not used to reading
mathematical proofs. However, we feel justified in concluding that the
program’s output is not _obviously_ written by a computer program, and that
was our objective.
## 4 Future work
In the long term, we would like to enlarge significantly the set of problems
that our program, or some new version of it, is capable of solving. To do
this, we will have to enable the program to handle certain kinds of deductions
that it currently handles either not at all or only in a rather rudimentary
way. In particular, an immediate target is to give the program the means to
deal with second-order quantification, which would allow it to solve easy
compactness problems, and also problems that require the construction of
‘obvious’ sequences.
At a more basic level, the program does not currently solve problems that
involve proof by contraposition or contradiction. It is not hard to add moves
that allow it to cope with a few problems of this kind, but it is trickier to
do so while not letting it apply those moves in inappropriate contexts. More
work is needed to understand what triggers the ‘contradiction move’ in human
mathematicians, but we expect to be able to add this facility in the near
future.
The program is also not as good as we would like at handling equality
substitutions. The situation here is similar: we can obviously add moves that
do such substitutions (and have done so in the current version of the
program), but it is more challenging to understand when humans make such
substitutions. It is also tricky to come up with a general understanding of
how they choose which out of two equal variables or complex terms to
eliminate. At its most general, the problem of how to handle equality is well
known to be hard, but our immediate aim would be a program that can handle the
easy cases of that problem competently and in a human way.
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* (44) Wos, L., Pereira, F., Hong, R., Boyer, R.S., Moore, J.S., Bledsoe, W.W., Henschen, L., Buchanan, B.G., Wrightson, G., Green, C.: An overview of automated reasoning and related fields. Journal of Automated Reasoning 1(1), 5–48 (1985)
|
arxiv-papers
| 2013-09-17T22:56:06 |
2024-09-04T02:49:51.048672
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "M. Ganesalingam and W. T. Gowers",
"submitter": "Timothy Gowers",
"url": "https://arxiv.org/abs/1309.4501"
}
|
1309.4602
|
S. Bhattacharya, P. Chalermsook, K. Mehlhorn, and A. Neumann
# New Approximability Results for the Robust $k$-Median Problem
Sayan Bhattacharya Max-Planck Institut für Informatik
Campus E1 4,
66123 Saarbrücken, Germany
{bsayan,parinya,mehlhorn,aneumann}@mpi-inf.mpg.de Parinya Chalermsook Max-
Planck Institut für Informatik
Campus E1 4,
66123 Saarbrücken, Germany
{bsayan,parinya,mehlhorn,aneumann}@mpi-inf.mpg.de Kurt Mehlhorn Max-Planck
Institut für Informatik
Campus E1 4,
66123 Saarbrücken, Germany
{bsayan,parinya,mehlhorn,aneumann}@mpi-inf.mpg.de Adrian Neumann Max-Planck
Institut für Informatik
Campus E1 4,
66123 Saarbrücken, Germany
{bsayan,parinya,mehlhorn,aneumann}@mpi-inf.mpg.de
###### Abstract.
We consider a robust variant of the classical $k$-median problem, introduced
by Anthony et al. [2]. In the _Robust $k$-Median problem_, we are given an
$n$-vertex metric space $(V,d)$ and $m$ client sets $\left\\{S_{i}\subseteq
V\right\\}_{i=1}^{m}$. The objective is to open a set $F\subseteq V$ of $k$
facilities such that the worst case connection cost over all client sets is
minimized; in other words, minimize $\max_{i}\sum_{v\in S_{i}}d(F,v)$. Anthony
et al. showed an $O(\log m)$ approximation algorithm for any metric and APX-
hardness even in the case of uniform metric. In this paper, we show that their
algorithm is nearly tight by providing $\Omega(\log m/\log\log m)$
approximation hardness, unless ${\sf NP}\subseteq\bigcap_{\delta>0}{\sf
DTIME}(2^{n^{\delta}})$. This hardness result holds even for uniform and line
metrics. To our knowledge, this is one of the rare cases in which a problem on
a line metric is hard to approximate to within logarithmic factor. We
complement the hardness result by an experimental evaluation of different
heuristics that shows that very simple heuristics achieve good approximations
for realistic classes of instances.
###### Key words and phrases:
Hardness of Approximation, Heuristics
###### 1991 Mathematics Subject Classification:
F.2.2 Nonnumerical Algorithms and Problems
## 1\. Introduction
In the classical $k$-median problem, we are given a set of clients located on
a metric space with distance function $d:V\times V\rightarrow\mathbb{R}$. The
goal is to open a set of facilities $F\subseteq V$, $|F|=k$, so as to minimize
the sum of the connection costs of the clients in $V$, i.e., their distances
from their nearest facilities in $F$. This is a central problem in
approximation algorithms, and quite naturally, it has received a large amount
of attention in the past two decades [6, 5, 7, 13, 12].
At SODA 2008 Anthony et al. [1, 2] introduced a generalization of the
$k$-median problem. In their setting, the set of clients that are to be
connected to some facility is not known in advance, and the goal is to perform
well in spite of this uncertainty about the future. In particular, they
formulated the problem as follows.
###### Definition 1.1 (Robust $k$-Median).
An instance of this problem is a triple $(V,{\mathcal{S}},d)$. This defines a
set of _locations_ $V$, a collection of $m$ sets of _clients_
${\mathcal{S}}=\\{S_{1},\ldots,S_{m}\\}$, where $S_{i}\subseteq V$ for all
$i\in\\{1,\ldots,m\\}$, and a metric distance function $d:V\times
V\rightarrow\mathbb{R}$. We have to open a set of $k$ facilities $F\subseteq
V$, $|F|=k$, and the goal is to minimize the cost of the most expensive set of
clients, i.e. minimize $\max_{i=1}^{m}\sum_{v\in S_{i}}d(v,F)$. Here, $d(v,F)$
denotes the minimum distance of the client $v$ from any location in $F$, i.e.
$d(v,F)=\min_{u\in F}d(u,v)$.
Note that the Robust $k$-Median problem is a natural generalization of the
classical $k$-median problem (where $m=1$). In addition, we can think of this
formulation as capturing some notion of fairness. To see this, simply
interpret each set $S_{i}$ as a community of clients who would pay $\sum_{v\in
S_{i}}d(v,F)$ for getting connected to some facility. Now the objective
ensures that no single community pays too much, while minimizing the cost.
Anthony et al. [2] gave an $O(\log m)$-approximation algorithm for this
problem, and a lower bound of $(2-\epsilon)$ for the best possible
approximation ratio by a reduction from Vertex Cover.
#### Our Results
We give nearly tight hardness of approximation results for the Robust
$k$-Median problem. We show that unless ${\sf NP}\subseteq\cap_{\delta>0}{\sf
DTIME}(2^{n^{\delta}})$, the problem admits no poly-time $o(\log m/\log\log
m)$-approximation, even on uniform and line metrics.
Our first hardness result is tight up to a constant factor, as a simple
rounding scheme gives a matching upper bound on uniform metrics (see Section
3.1). Our second, and rather surprising, result shows that “Robust $k$-Median”
is a rare problem with super-constant hardness of approximation even on line
metrics, in sharp contrast to most other geometric optimization problems which
admit polynomial time approximation schemes, e.g. [3, 11].
In Section 5 we investigate the performance of some heuristics. Already a very
simple greedy strategy provides reasonably good performance on a realistic
class of instances. We use an LP relaxation of the problem as a lower bound.
#### Our Techniques
First, we note that the Robust $k$-Median problem on uniform metrics is
equivalent to the following variant of the set cover problem: Given a set $U$
of ground elements, a collection of sets ${\mathcal{X}}=\left\\{X\subseteq
U\right\\}$, and an integer $t\leq|{\mathcal{X}}|$, our goal is to select $t$
sets from ${\mathcal{X}}$ in order to minimize the number of times an element
from $U$ is hit (see Lemma 3.2). We call this problem Minimum Congestion Set
Packing (MCSP). This characterization allows us to focus on proving the
hardness of MCSP, and to employ the tools developed for the set cover problem.
We now revisit the reduction used in proving the hardness of the set cover
problem by Feige [8], building on the framework of Lund and Yannakakis [14],
and discuss how our approach differs from theirs. Intuitively, they compose
the Label Cover instance with a set system that has some desirable properties.
Informally speaking, in the Label Cover problem, we are given a graph where
each vertex $v$ can be assigned a label from a set $L$, and each edge $e$ is
equipped with a constraint $\Pi_{e}\subseteq L\times L$ specifying the
accepting pairs of labels for $e$. Our goal is to find a labeling of vertices
that maximizes the number of accepting edges. This problem is known to be hard
to approximate to within a factor of $2^{\log^{1-\epsilon}|E|}$ [4, 15], where
$|E|$ is the number of edges. Thus, if we manage to reduce Label Cover to
MCSP, we would hopefully obtain a large hardness of approximation factor for
MCSP as well.
From the Label Cover instance, [14] creates an instance of Set Cover by having
sets of the form $S(v,\ell)$ for each vertex $v$ and each label $\ell\in L$.
Intuitively the set $S(v,\ell)$ means choosing label $\ell$ for vertex $v$ in
the label cover instance. Now, if we assume that the solution is well behaved,
in the sense that for each vertex $v$, only one set of the form $S(v,\ell)$ is
chosen in the solution, we would be immediately done (because each set indeed
corresponds to label assignment). However, a solution need not have this form,
e.g. choosing sets $S(v,\ell)$ and $S(v,\ell^{\prime})$ would translate to
having two labels $\ell,\ell^{\prime}$ for the label cover instance. To
prevent an ill-behaved solution, “partition systems” were introduced and used
in both [14] and [8]. Feige considers the hypergraph version of Label Cover to
obtain a sharper hardness result of $\ln n-O(\ln\ln n)$ instead of
$\frac{1}{4}\ln n$ in [14]; here $n$ denotes the size of the universe.
Now we highlight how our reduction is different from theirs. The high level
idea of our reduction is the same, i.e. we have sets of the form $S(v,\ell)$
that represent assigning label $\ell$ to vertex $v$. However, we need a
different partition system and a totally different analysis. Moreover, while a
reduction from standard Label Cover gives nearly tight $O(\log n)$ hardness
for Set Cover, it can (at best) only give the hardness of $2-\epsilon$ for
MCSP. To prove our results, we do need a reduction from the Hypergraph Label
Cover problem. This suggests another natural distinction between MCSP and Set
Cover.
Finally, to obtain the hardness of the Robust $k$-Median problem on the line
metric, we embed the instance created from the MCSP reduction onto the line
such that the values of optimal solutions are preserved. This way we get the
same hardness gap for line metrics.
## 2\. Preliminaries
We will show that the Robust $k$-Median problem is $\Omega(\log m/\log\log m)$
hard to approximate, even for the special cases of _uniform metrics_ (see
Section 3) and _line metrics_ (see Section 4). Recall that $d$ is a uniform
metric iff we have $d(u,v)\in\\{0,1\\}$ for all locations $u,v\in V$. Further,
$d$ is a line metric iff the locations in $V$ can be embedded into a line in
such a way that $d(u,v)$ equals the Euclidean distance between $u$ and $v$,
for all $u,v\in V$. Throughout this paper, we will denote any set of the form
$\\{1,2,\ldots,i\\}$ by $[i]$. Our hardness results will rely on a reduction
from the _$r$ -Hypergraph Label Cover_ problem, which is defined as follows.
###### Definition 2.1 ($r$-Hypergraph Label Cover).
An instance of this problem is a triple $(G,\pi,r)$, where
$G=(\mathcal{V},\mathcal{E})$ is a $r$-partite hypergraph with vertex set
${\mathcal{V}}=\bigcup_{j=1}^{r}{\mathcal{V}}_{j}$ and edge set
${\mathcal{E}}$. Each edge $h\in\mathcal{E}$ contains one vertex from each
part of ${\mathcal{V}}$, i.e. $|h\cap{\mathcal{V}}_{j}|=1$ for all $j\in[r]$.
Every set ${\mathcal{V}}_{j}$ has an associated set of _labels_ $L_{j}$.
Further, for all $h\in\mathcal{E}$ and $j\in[r]$, there is a mapping
$\pi_{h}^{j}:L_{j}\rightarrow C$ that projects the labels from $L_{j}$ to a
common set of _colors_ $C$.
The problem is to assign to every vertex $v\in{\mathcal{V}}_{j}$ some label
$\sigma(v)\in L_{j}$. We say that an edge $h=(v_{1},\ldots,v_{r})$, where
$v_{j}\in{\mathcal{V}}_{j}$ for all $j\in[r]$, is _strongly satisfied_ under
$\sigma$ iff the labels of all its vertices are mapped to the same element in
$C$, i.e.
$\pi_{h}^{j}(\sigma(v_{j}))=\pi_{h}^{j^{\prime}}(\sigma(v_{j^{\prime}}))$ for
all $j,j^{\prime}\in[r]$. In contrast, we say that the edge is _weakly
satisfied_ iff there exists some pair of vertices in $h$ whose labels are
mapped to the same element in $C$, i.e.
$\pi^{j}_{h}(\sigma(v_{j}))=\pi^{j^{\prime}}_{h}(\sigma(v_{j^{\prime}}))$ for
some $j,j^{\prime}\in[r]$, $j\neq j^{\prime}$.
For ease of exposition, we will often abuse the notation and denote by $j(v)$
the part of ${\mathcal{V}}$ to which a vertex $v$ belongs, i.e. if
$v\in{\mathcal{V}}_{j}$ for some $j\in[r]$, then we set $j(v)\leftarrow j$.
The next theorem will be crucial in deriving our hardness result. The proof of
this theorem follows from Feige’s $r$-Prover system [8] (see Appendix A).
###### Theorem 2.2.
Let $r\in\mathbb{N}$ be a parameter. There is a polynomial time reduction from
$n$-variable 3-SAT to $r$-Hypergraph Label Cover with the following
properties:
* •
(Yes-Instance) If the formula is satisfiable, then there is a labeling that
strongly satisfies every edge in $G$.
* •
(No-Instance) If the formula is not satisfiable, then every labeling weakly
satisfies at most a $2^{-\gamma r}$ fraction of the edges in $G$, for some
universal constant $\gamma$.
* •
The number of vertices in the graph is $|{\mathcal{V}}|=n^{O(r)}$ and the
number of edges is $|\mathcal{E}|=n^{O(r)}$. The sizes of the label sets are
$\left|L_{j}\right|=2^{O(r)}$ for all $j\in[r]$, and
$\left|C\right|=2^{O(r)}$. Further, we have
$|{\mathcal{V}}_{j}|=|{\mathcal{V}}_{j^{\prime}}|$ for all
$j,j^{\prime}\in[r]$, and each vertex $v\in{\mathcal{V}}$ has the same degree
$r|\mathcal{E}|/|{\mathcal{V}}|$.
We use a _partition system_ that is motivated by the hardness proof of the Set
Cover problem [8]. However, we deal with a different problem, and our
construction is also different.
###### Definition 2.3 (Partition System).
Let $r\in\mathbb{N}$ and let $C$ be any finite set. An $(r,C)$-partition
system is a pair $(Z,\left\\{p_{c}\right\\}_{c\in C})$, where $Z$ is an
arbitrary (ground) set, and for each $c\in C$, $p_{c}$ is a partition of $Z$
into $r$ subsets, such that the following properties hold.
* •
(Partition) For all $c\in C$, $p_{c}=\left(A_{c}^{1},\ldots,A_{c}^{r}\right)$
is a partition of $Z$, that is $\bigcup_{j=1}^{r}A_{c}^{j}=Z$, and
$A_{c}^{j^{\prime}}\cap A_{c}^{j}=\emptyset$ for all $j,j^{\prime}\in[r],j\neq
j^{\prime}$.
* •
($r$-intersecting) For any $r$ _distinct_ indices $c_{1},\ldots,c_{r}\in C$
and _not-necessarily distinct_ indices $j_{1},\ldots,j_{r}\in[r]$, we have
that $\bigcap_{i=1}^{r}A_{c_{i}}^{j_{i}}\neq\emptyset$. In particular,
$A_{c}^{j}\not=\emptyset$ for all $c$ and $j$.
In order to achieve a good lower bound on the approximation factor, we need
partition systems with _small_ ground sets. The most obvious way to build a
partition system is to form an $r$-hypercube: Let $Z=[r]^{|C|}$, and for each
$c\in C$ and $j\in[r]$, let $A_{c}^{j}$ be the set of all elements in $Z$
whose $c$-th component is $j$. It can easily be verified that this is an
$(r,C)$-partition system with $|Z|=r^{\left|C\right|}$. With this
construction, however, we would only get a hardness of $\Omega(\log\log m)$
for our problem. The following lemma shows that it is possible to construct an
$(r,C)$-partition system probabilistically with
$|Z|=r^{O(r)}\log\left|C\right|$.
###### Lemma 2.4.
There is an $(r,C)$-partition system with $|Z|=r^{O(r)}\log\left|C\right|$
elements. Further, such a partition system can be constructed efficiently with
high probability.
###### Proof 2.5.
Let $Z$ be any set of $r^{O(r)}\log\left|C\right|$ elements. We build a
partition system $(Z,\\{p_{c}\\}_{c\in C})$ as described in Algorithm 1.
input : A ground set $Z$, a parameters $r\in\mathbb{N}$, and a set $C$.
foreach _$c\in C$_ do
/* Construct the partition $p_{c}=(A_{c}^{1},\ldots,A_{c}^{r})$ */
Initialize $A_{c}^{j}$ to the empty set for all $j\in[r]$
foreach _ground element $e\in Z$_ do
Pick an index $j\in[r]$ independently and uniformly at random and add $e$ to
$A_{c}^{j}$
Algorithm 1 A randomized algorithm for constructing an $(r,C)$-partition
system.
In Algorithm 1, by construction each $p_{c}$ is a partition of $Z$, i.e. the
first property stated in Definition 2.3 is satisfied. We bound the probability
that the second property is violated.
Fix any choice of $r$ _distinct_ indices $c_{1},\ldots,c_{r}\in C$ and _not
necessarily distinct_ indices $j_{1},\ldots,j_{r}\in[r]$. We say that a _bad
event_ occurs when the intersection of the corresponding sets is empty, i.e.
$\bigcap_{i=1}^{r}A_{c_{i}}^{j_{i}}=\emptyset$. To upper bound the probability
of a bad event, we focus on events of the form $E_{e,i}$ – this occurs when an
element $e\in Z$ is included in a set $A_{c_{i}}^{j_{i}}$. Since the indices
$c_{1}\ldots c_{r}$ are distinct, it follows that the events $\\{E_{e,i}\\}$
are mutually independent. Furthermore, note that we have $\Pr[E_{e,i}]=1/r$
for all $e\in Z,i\in[r]$. Hence, the probability that an element $e\in Z$ does
not belong to the intersection $\bigcap_{i=1}^{r}A_{c_{i}}^{j_{i}}$ is given
by $1-\Pr[\bigcap_{i=1}^{r}E_{e,i}]=1-1/r^{r}$. Accordingly, the probability
that no element $e\in Z$ belongs to the intersection, which defines the bad
event, is equal to $(1-1/r^{r})^{|Z|}$.
Now, the total number of choices for $r$ distinct indices $c_{1},\ldots,c_{r}$
and $r$ not-necessarily distinct indices $j_{1},\ldots,j_{r}$ is equal to
$\binom{\left|C\right|}{r}\cdot r^{r}$. Hence, taking a union-bound over all
possible bad events, we see that the second property stated in Definition 2.3
is violated with probability at most $\binom{\left|C\right|}{r}\cdot
r^{r}\cdot(1-r^{r})^{|Z|}\leq(\left|C\right|r)^{r}\cdot\exp(-|Z|/r^{r})$. If
we set $|Z|=d\cdot r^{d\cdot r}\log\left|C\right|$ with sufficiently large
constant $d$, then it is easy to see that the second constraint in Definition
2.3 is satisfied with high probability.
## 3\. Hardness of Robust k-Median on Uniform Metrics
First, we define a problem called _Minimum Congestion Set Packing_ (MCSP), and
then show a reduction from MCSP to Robust $k$-Median on uniform metrics. In
Section 3.2, we will then show that MCSP is hard to approximate by reducing
Hypergraph Label Cover to MCSP.
###### Definition 3.1.
[Minimum Congestion Set Packing (MCSP)] An instance of this problem is a
triple $(U,{\mathcal{X}},t)$, where $U$ is a universe of $m$ elements, i.e.
$|U|=m$, ${\mathcal{X}}$ is a collection of sets
${\mathcal{X}}=\left\\{X\subseteq U\right\\}$ such that
$\bigcup_{X\in{\mathcal{X}}}X=U$, and $t\in\mathbb{N}$ and
$t\leq\left|{\mathcal{X}}\right|$. The objective is to find a collection
${\mathcal{X}}^{\prime}\subseteq{\mathcal{X}}$ of size $t$ that minimizes
$\text{{Cong}}({\mathcal{X}}^{\prime})=\max_{e\in
U}\text{{Cong}}(e,{\mathcal{X}}^{\prime})$. Here,
$\text{{Cong}}({\mathcal{X}}^{\prime})$ refers to the _congestion_ of the
solution ${\mathcal{X}}^{\prime}$, and
$\text{{Cong}}(e,{\mathcal{X}}^{\prime})=|\\{X\in{\mathcal{X}}^{\prime}:e\in
X\\}|$ is the congestion of the element $e\in U$ under the solution
${\mathcal{X}}^{\prime}$.
###### Lemma 3.2.
Given any MCSP instance $(U,{\mathcal{X}},t)$, we can construct a Robust
$k$-Median instance $(V,{\mathcal{S}},d)$ with the same objective value in
$poly(|U|,|{\mathcal{X}}|)$ time, such that $|U|=|{\mathcal{S}}|$,
$|{\mathcal{X}}|=|V|$, $d$ is a uniform metric, and $k=|V|-t$.
###### Proof 3.3.
We construct the Robust $k$-Median instance $(V,{\mathcal{S}},d)$ as follows.
For every $e\in U$ we create a set of clients $S(e)$, and for each
$X\in{\mathcal{X}}$ we create a location $v(X)$. Thus, we get
$V=\\{v(X)\,:\,X\in{\mathcal{X}}\\}$, and ${\mathcal{S}}=\\{S(e)\,:\,e\in
U\\}$. We place the clients in $S(e)$ at the locations of the sets that
contain $e$, i.e. $S(e)=\\{v(X)\,:\,X\in{\mathcal{X}},e\in X\\}$ for all $e\in
U$. The distance is defined as $d(u,v)=1$ for all $u,v\in V,u\neq v$, and
$d(v,v)=0$. Finally, we set $k\leftarrow|V|-t$.
Now, it is easy to verify that the Robust $k$-Median instance
$(V,{\mathcal{S}},d)$ has a solution with objective $\rho$ iff the
corresponding MCSP instance $(U,{\mathcal{X}},t)$ has a solution with
objective $\rho$. The intuition is that a location $v(X)\in V$ is _not_
included in the solution $F$ to the Robust $k$-Median instance iff the
corresponding set $X$ is included in the solution ${\mathcal{X}}^{\prime}$ to
the MCSP instance. Indeed, let $F$ be any subset of ${\mathcal{X}}$ of size
$k$ (= the set of open facilities) and let
${\mathcal{X}}^{\prime}={\mathcal{X}}-F$. Further, let
$[X\in{\mathcal{X}}^{\prime}]$ be an indicator variable that is set to $1$ iff
$X\in{\mathcal{X}}^{\prime}$. Then
$\displaystyle\text{{Cong}}({\mathcal{X}}^{\prime})$ $\displaystyle=\max_{e\in
U}\text{{Cong}}(e,{\mathcal{X}}^{\prime})=\max_{e\in U}\sum_{X;e\in
X}[X\in{\mathcal{X}}^{\prime}]$ $\displaystyle=\max_{e\in U}\sum_{X;e\in
X}\min_{Y\in F}d(X,Y)=\max_{S(e)\in{\mathcal{S}}}\sum_{v(X)\in
S(e)}d(v(X),F).$
We devote the rest of Section 3 to the MCSP problem and show that it is
$\Omega(\log|U|/\log\log|U|)$ hard to approximate. This, in turn, will imply a
$\Omega(\log|{\mathcal{S}}|/\log\log|{\mathcal{S}}|)$ hardness of
approximation for Robust $k$-Median on uniform metrics. We will prove the
hardness result via a reduction from Hypergraph Label Cover.
### 3.1. Integrality Gap
Before proceeding to the hardness result, we show that a natural LP relaxation
for the MCSP problem [2] has an integrality gap of $\Omega(\log m/\log\log
m)$, where $m=|U|$ is the size of the universe of elements. In the LP, we have
a variable $y(X)$ indicating that the set $X\in{\mathcal{X}}$ is chosen, and a
variable $z$ which represents the maximum congestion among the elements.
$\displaystyle\min\quad$ $\displaystyle z$ s.t.
$\displaystyle\sum_{\mathclap{X\in{\mathcal{X}}:e\in X}}y(X)\leq z\mbox{ for
all $e\in U$}$ $\displaystyle\sum_{\mathclap{X\in{\mathcal{X}}}}y(X)=t$
The Instance: Now, we construct a bad integrality gap instance
$(U,{\mathcal{X}},t)$. Let $d$ be the intended integrality gap, let
$\eta=d^{2}$, and let
$U=\left\\{I\,:\,I\subseteq[\eta],\left|I\right|=d\right\\}$ be all subsets of
$[\eta]$ of size $d$. The collection ${\mathcal{X}}$ consists of $\eta$ sets
$X_{1},\ldots,X_{\eta}$, where $X_{i}=\left\\{I\,:\,I\in U\text{ and }i\in
I\right\\}$. Note that the universe $U$ consists of $|U|=m=\binom{\eta}{d}$
elements, and each element $I$ is contained in exactly $d$ sets, namely $I\in
X_{i}$ if and only if $i\in I$. Finally, we set $t\leftarrow\eta/d$.
Analysis: The fractional solution simply assigns a value of $1/d$ to each
variable $y(X_{i})$; this ensures that the total (fractional) number of sets
selected is $\eta/d=t$. Furthermore, each element is contained (fractionally)
in exactly one set, so the fractional solution has cost one. Any integral
solution must choose $\eta/d=d$ sets, say $X_{i_{1}}\ldots X_{i_{d}}$. Then
$I=\left\\{i_{1},\ldots,i_{d}\right\\}\in X_{i_{\lambda}}$ for all
$\lambda\in[d]$ and hence the congestion of $I$ is $d$, and this also means
that any integral solution has cost at least $d$. Finally, since
$|U|=m\leq\eta^{d}\leq(d^{2})^{d}$, we have $d=\Omega(\log m/\log\log m)$.
Tightness of the result: The bound on the hardness and integrality gap is
tight for the uniform metric case, as there is a simple $O(\log m/\log\log
m)$-approximation algorithm. Pick each set $X$ with probability equal to
$\min(1,2y(X))$. The expected congestion is $2z$ for each element. By
Chernoff’s bound [10], an element is covered by no more than $z\cdot O(\log
m/\log\log m)$ sets with high probability. A similar algorithm gives the same
approximation guarantee for the Robust $k$-Median problem on uniform metrics.
### 3.2. Reduction from r-Hypergraph Label Cover to Minimum Congestion Set
Packing
The input is an instance $\left(G,\pi,r\right)$ of the $r$-Hypergraph Label
Cover problem (see Definition 2.1). From this we construct the following
instance $(U,{\mathcal{X}},t)$ of the MCSP problem (see Definition 3.1).
* •
First, we define the universe $U$ as a union of disjoint sets. For each edge
$h\in{\mathcal{E}}$ in the hypergraph we have a set $U_{h}$. All these sets
have the same size $m^{*}$ and are pairwise disjoint, i.e. $U_{h}\cap
U_{h^{\prime}}=\emptyset$ for all $h,h^{\prime}\in{\mathcal{E}}$,
$h^{\prime}\neq h$. The universe $U$ is then the union of these sets
$U=\bigcup_{h\in{\mathcal{E}}}U_{h}$. Since the $U_{h}$ are mutually disjoint,
we have $m=|U|=|{\mathcal{E}}|\cdot m^{*}$. Recall that $C$ is the target set
of $\pi$. Each set $U_{h}$ is the ground set of an $(r,C)$-partition system
(see Definition 2.3) as given by Lemma 2.4. In particular we have
$m^{*}=r^{O(r)}\log\left|C\right|$. We denote the $r$-partitions associated
with $U_{h}$ by $\\{p_{c}(h)\\}_{c\in C}$, where
$p_{c}(h)=\left(A^{1}_{c}(h),\ldots,A^{r}_{c}(h)\right)$.
* •
Second, we construct the collection of sets ${\mathcal{X}}$ as follows. For
each $j\in[r]$, $v\in{\mathcal{V}}_{j}$ and $\ell\in L_{j}$, ${\mathcal{X}}$
contains the set $X(v,\ell)$, where $X(v,\ell)=\bigcup_{h:v\in
h}A^{j}_{\pi^{j}_{h}(\ell)}(h)$. In other words, $X(v,\ell)\cap U_{h}$ is
empty if $v\not\in h$ and is equal to $A^{j}_{\pi^{j}_{h}(\ell)}(h)$ if $v\in
h$. Intuitively, choosing the set $X(v,\ell)$ corresponds to assigning label
$\ell$ to the vertex $v$.
* •
Third, we define $t\leftarrow|{\mathcal{V}}|$. Intuitively, this means that
each vertex in ${\mathcal{V}}$ gets one label.
We assume for the sequel that the $r$-Hypergraph Label Cover instance is
chosen according to Theorem 2.2. We assume that the parameter $r$ satisfies
$r^{7}2^{-\gamma r}<1$. In the proof of the main theorem, we will fix $r$ to a
specific value.
### 3.3. Analysis
We show that the reduction from Hypergraph Label Cover to MCSP satisfies two
properties. In Lemma 3.4, we show that for a Yes-Instance (see Theorem 2.2),
the corresponding MCSP instance admits a solution with congestion one. Second,
in case of a No-Instance, we show in Lemma 3.10 that every solution to the
corresponding MCSP instance has congestion at least $r$.
###### Lemma 3.4 (Yes-Instance).
If the Hypergraph Label Cover instance $(G,\pi,r)$ admits a labeling that
strongly satisfies every edge, then the MCSP instance $(U,{\mathcal{X}},t)$
constructed in Section 3.2 admits a solution where the congestion of every
element in $U$ is exactly one.
###### Proof 3.5.
Suppose that there is a labeling $\sigma$ that strongly satisfies every edge
$h\in{\mathcal{E}}$. We will show how to pick $t=|{\mathcal{V}}|$ sets from
${\mathcal{X}}$ such that each element in $U$ is contained in exactly one set.
This implies that the maximum congestion is one. For each $j\in[r]$ and each
vertex $v\in{\mathcal{V}}_{j}$, we choose the set $X(v,\sigma(v))$. Thus, the
total number of sets chosen is exactly $|{\mathcal{V}}|$.
To see that the congestion is indeed one, we concentrate on the elements in
$U_{h}$, where $h=(v_{1},\ldots,v_{r})$, $v_{j}\in{\mathcal{V}}_{j}$ for all
$j\in[r]$, is one of the edges in ${\mathcal{E}}$. The picked sets that
intersect $U_{h}$ are $X(v_{j},\sigma(v_{j}))$, where $j\in[r]$. Since $h$ is
strongly satisfied, $\pi_{h}$ maps all labels of the vertices in $h$ to a
common $c\in C$, that is $\pi^{j}_{h}(\sigma(v_{j}))=c$ for all $j\in[r]$.
Thus $U_{h}\cap X(v_{j},\sigma(v_{j}))=A^{j}_{c}(h)$. By the definition of a
partition system (see Definition 2.3), the sets $A^{1}_{c}(h)\ldots
A^{r}_{c}(h)$ partition the elements in $U_{h}$. This completes the proof.
Now, we turn to the proof of Lemma 3.10. Towards this end, we fix a collection
${\mathcal{X}}^{\prime}\subseteq{\mathcal{X}}$ of size $t$ and show that some
element in $U$ has congestion at least $r$ under ${\mathcal{X}}^{\prime}$. The
intuition being that many edges in $G=({\mathcal{V}},{\mathcal{E}})$ are not
even weakly satisfied, and the elements in $U$ corresponding to those edges
incur large congestion. Recall that for a $v\in{\mathcal{V}}$, we define
$j(v)\in\mathbb{N}$ to be such that $v\in{\mathcal{V}}_{j(v)}$.
###### Claim 1.
For $v\in{\mathcal{V}}$, let ${\mathcal{L}}_{v}=\left\\{\ell\in
L_{j(v)}\,:\,X(v,\ell)\in{\mathcal{X}}^{\prime}\right\\}$. For
$h\in{\mathcal{E}}$, let
$\Lambda_{h}=\left\\{X(v,\ell)\in{\mathcal{X}}^{\prime}\,:\,v\in h\right\\}$
and $\lambda(h)=\left|\Lambda_{h}\right|$. If the solution
${\mathcal{X}}^{\prime}$ has congestion less than $r$ then
$\left|{\mathcal{L}}_{v}\right|<r^{2}$ and $\left|\Lambda_{h}\right|<r^{3}$.
###### Proof 3.6.
Since $\Lambda_{h}=\bigcup_{v\in h}{\mathcal{L}}_{v}$, it suffices to prove
$\left|{\mathcal{L}}_{v}\right|<r^{2}$ for all $v$. Assume otherwise, i.e.,
$\left|{\mathcal{L}}_{v}\right|\geq r^{2}$ for some $v\in{\mathcal{V}}_{j}$,
$j\in[r]$. Let $h$ be any hyper-edge with $v\in h$. Consider the images of the
labels in ${\mathcal{L}}_{v}$ under the projection $\pi^{j}_{h}$. Either we
have at least $r$ distinct images or at least $r$ elements in $L_{v}$ are
mapped to the same element of $C$.
In the former case, we have $r$ pairwise distinct labels $\ell_{1}$ to
$\ell_{r}$ in ${\mathcal{L}}_{v}$ and $r$ pairwise distinct labels $c_{1}$ to
$c_{r}$ in $C$ such that $\pi^{j}_{h}(\ell_{i})=c_{i}$ for $i\in[r]$. The set
$X(v,\ell_{i})$ contains $A^{j}_{c_{i}}(h)$ and
$\bigcap_{i\in[r]}A^{j}_{c_{i}}(h)\not=\emptyset$ by property (2) of partition
systems (see Definition 2.3). Thus some element has congestion at least $r$.
In the latter case, we have $r$ pairwise distinct labels $\ell_{1}$ to
$\ell_{r}$ in ${\mathcal{L}}_{v}$ and a label $c$ in $C$ such that
$\pi^{j}_{h}(\ell_{i})=c$ for $i\in[r]$. The set $X(v,\ell_{i})$ contains
$A^{j}_{c}(h)$ and hence every element in this non-empty set (property (2) of
partition systems) has congestion at least $r$.
###### Definition 3.7 (Colliding Edge).
We say that an edge $h\in{\mathcal{E}}$ is _colliding_ iff there are sets
$X(v,\ell),X(v^{\prime},\ell^{\prime})\in{\mathcal{X}}^{\prime}$ with
$v,v^{\prime}\in h$, $v\not=v^{\prime}$, and
$\pi^{j(v)}_{h}(\ell)=\pi^{j(v^{\prime})}_{h}(\ell^{\prime})$.
###### Claim 2.
Suppose that the solution ${\mathcal{X}}^{\prime}$ has congestion less than
$r$, and more than a $r^{4}2^{-\gamma r}$ fraction of the edges in
${\mathcal{E}}$ are colliding. Then there is a labeling $\sigma$ for $G$ that
weakly satisfies at least a $2^{-\gamma r}$ fraction of the edges in
${\mathcal{E}}$.
###### Proof 3.8.
For each $v\in{\mathcal{V}}$, we define the label set
${\mathcal{L}}_{v}=\\{\ell\in
L_{j(v)}\,:\,X(v,\ell)\in{\mathcal{X}}^{\prime}\\}$. Then
$\left|{\mathcal{L}}_{v}\right|<r^{2}$ by Claim 1. We construct a labeling
function $\sigma$ using Algorithm 2.
foreach _vertex $v\in{\mathcal{V}}$_ do
if _${\mathcal{L}}_{v}\neq\emptyset$_ then
Pick a color $\sigma(v)$ uniformly and independently at random from
${\mathcal{L}}_{v}$
else
Pick an arbitrary color $\sigma(v)$ from $L_{j(v)}$
Algorithm 2 An algorithm for constructing a labeling function.
Now we bound the expected fraction of weakly satisfied edges under $\sigma$
from below. Take any colliding edge $h\in{\mathcal{E}}$. This means that there
are vertices $v\in{\mathcal{V}}_{j}$,
$v^{\prime}\in{\mathcal{V}}_{j^{\prime}}$ with $j\neq j^{\prime}$, and colors
$\ell\in{\mathcal{L}}_{v}$, $\ell^{\prime}\in{\mathcal{L}}_{v^{\prime}}$ such
that $v,v^{\prime}\in h$ and
$\pi^{j}_{h}(\ell)=\pi^{j^{\prime}}_{h}(\ell^{\prime})$. By Claim 1,
$|{\mathcal{L}}_{v}|$ and $|{\mathcal{L}}_{v^{\prime}}|$ are both at most
$r^{2}$. Since the colors $\sigma(v)$ and $\sigma(v^{\prime})$ are chosen
uniformly and independently at random from their respective palettes
${\mathcal{L}}_{v}$ and ${\mathcal{L}}_{v^{\prime}}$, we have
$\Pr[\sigma(v)=\ell\text{ and }\sigma(v^{\prime})=\ell^{\prime}]\geq 1/r^{4}$.
In other words, every colliding edge is weakly satisfied with probability at
least $1/r^{4}$. Since more than a $r^{4}2^{-\gamma r}$ fraction of the edges
in ${\mathcal{E}}$ are colliding, from linearity of expectation we infer that
the expected fraction of edges weakly satisfied by $\sigma$ is at least
$2^{-\gamma r}$.
###### Claim 3.
Let $\Lambda_{h}=\left\\{X(v,\ell)\in{\mathcal{X}}^{\prime}\,:\,v\in
h\right\\}$, and $\lambda(h)=|\Lambda_{h}|$. We have
$\sum_{h\in{\mathcal{E}}}\lambda(h)=r|{\mathcal{E}}|$.
###### Proof 3.9.
This is a simple counting argument. Consider a bipartite graph $H$ with vertex
set $A\dot{\cup}B$, where each vertex in $A$ represents a set $X(v,\ell)$, and
each vertex in $B$ represents an edge $h\in{\mathcal{E}}$. There is an edge
between two vertices iff the set $X(v,\ell)$ contains some element in $U_{h}$.
The quantity $\sum_{h\in{\mathcal{E}}}\lambda(h)$ counts the number of edges
in $H$ where one endpoint is included in the solution
${\mathcal{X}}^{\prime}$. Since ${\mathcal{X}}^{\prime}$ picks
$t=|{\mathcal{V}}|$ sets and each set has degree
$r|{\mathcal{E}}|/|{\mathcal{V}}|$ in the $H$ (see Theorem 2.2), the total
number of edges that are chosen is exactly
$|{\mathcal{V}}|\times\left(r|{\mathcal{E}}|/|{\mathcal{V}}|\right)=r|{\mathcal{E}}|$.
Let ${\mathcal{E}}^{\prime}\subseteq{\mathcal{E}}$ denote the set of colliding
edges, and define
${\mathcal{E}}^{\prime\prime}={\mathcal{E}}-{\mathcal{E}}^{\prime}$. Suppose
that we are dealing with a No-Instance (see Theorem 2.2), i.e. the solution
${\mathcal{X}}^{\prime}$ has congestion less than $r$ and every labeling
weakly satisfies at most a $2^{-\gamma r}$ fraction of the edges in
${\mathcal{E}}$. Then $\lambda(h)\leq r^{3}$ for all $h\in{\mathcal{E}}$ by
Claim 1, and no more than $r^{4}2^{-\gamma r}|{\mathcal{E}}|$ edges are
colliding, i.e. $\left|{\mathcal{E}}^{\prime}\right|\leq r^{4}2^{-\gamma
r}\left|{\mathcal{E}}\right|$, by Claim 2. Using these facts we conclude that
$\sum_{h\in{\mathcal{E}}^{\prime}}\lambda(h)\leq r^{7}2^{-\gamma
r}|{\mathcal{E}}|<|{\mathcal{E}}|$, as by assumption $r^{7}2^{-\gamma r}<1$.
Now, applying Claim 3, we get
$\sum_{h\in{\mathcal{E}}^{\prime\prime}}\lambda(h)=r|{\mathcal{E}}|-\sum_{h\in{\mathcal{E}}^{\prime}}\lambda(h)>(r-1)|{\mathcal{E}}|$.
In particular, there is an edge $h\in{\mathcal{E}}^{\prime\prime}$ with
$\lambda(h)\geq r$.
Recall that $\Lambda_{h}=\left\\{X(v,\ell)\in{\mathcal{X}}^{\prime}\,:\,v\in
h\right\\}$ are the sets in ${\mathcal{X}}^{\prime}$ that intersect $U_{h}$
and note that $|\Lambda_{h}|=\lambda(h)\geq r$. Let
${\mathcal{X}}^{*}\subseteq\Lambda_{h}$ be a _maximal_ collection of sets with
the following property: For every two distinct sets
$X(v,\ell),X(v^{\prime},\ell^{\prime})\in{\mathcal{X}}^{*}$ we have
$\pi^{j(v)}_{h}(\ell)\neq\pi^{j(v^{\prime})}_{h}(\ell^{\prime})$. Hence, from
the definition of a partition system (see Definition 2.3), it follows that the
intersection of the sets in ${\mathcal{X}}^{*}$ and the set $U_{h}$ is
nonempty.
Now, consider any set $X(v,\ell)\in\Lambda_{h}-{\mathcal{X}}^{*}$. Since the
collection ${\mathcal{X}}^{*}$ is maximal, there must be at least one set
$X(v^{\prime},\ell^{\prime})$ in ${\mathcal{X}}^{*}$ with
$\pi^{j(v)}_{h}(\ell)=\pi^{j(v^{\prime})}_{h}(\ell^{\prime})$. Since $h$ is
not colliding, we must have $j(v)=j(v^{\prime})$. Consequently we get
$X(v,\ell)\cap U_{h}=X(v^{\prime},\ell^{\prime})\cap U_{h}$. In other words,
for every set $X\in\Lambda_{h}-{\mathcal{X}}^{*}$, there is some set
$X^{\prime}\in{\mathcal{X}}^{*}$ where $X\cap U_{h}=X^{\prime}\cap U_{h}$.
Thus,
$U_{h}\cap(\bigcap_{X\in\Lambda_{h}}X)=U_{h}\cap(\bigcap_{X\in{\mathcal{X}}^{*}}X)\neq\emptyset$.
Every element in the intersection of the sets in $\Lambda_{h}$ and $U_{h}$
will have congestion $|\Lambda_{h}|\geq r$. This leads to the following lemma.
###### Lemma 3.10 (No-Instance).
Suppose that every labeling weakly satisfies at most a $2^{-\gamma r}$
fraction of the edges in the hypergragph label cover instance $(G,\pi,r)$, for
some universal constant $\gamma$ and that $r^{7}2^{-\gamma r}<1$. Then the
congestion incurred by every solution to the MCSP instance
$(U,{\mathcal{X}},t)$ constructed in Section 3.2 is at least $r$.
We are now ready to prove the main theorem of this section.
###### Theorem 3.11.
The Robust $k$-Median problem $(V,{\mathcal{S}},d)$ is $\Omega(\log m/\log\log
m)$ hard to approximate on uniform metrics, where $m=|{\mathcal{S}}|$, unless
$\textsf{NP}\subseteq\bigcap_{\delta>0}\textsf{DTIME}(2^{n^{\delta}})$.
###### Proof 3.12.
Assume that there is a polynomial time algorithm for the Robust $k$-Median
problem that guarantees an approximation ratio in
$o(\log\left|{\mathcal{S}}\right|/\log\log\left|{\mathcal{S}}\right|)$. Then,
by Lemma 3.2, there is an approximation algorithm for the Minimum Congestion
Set Packing problem with approximation guarantee
$o(\log\left|U\right|/\log\log\left|U\right|)$.
Let $\delta>0$ be arbitrary and set $r=\lfloor n^{\delta}\rfloor$, where $n$
is the number of variables in the $3$-SAT instance (see Theorem 2.2). Then
$r^{7}2^{-\gamma r}<1$ for all sufficiently large $n$. We first bound the size
of the MCSP instance $(U,{\mathcal{X}},t)$ constructed in Section 3.2. By
Lemma 2.4, the size of an $(r,C)$-partition system is
$|Z|=r^{O(r)}\log\left|C\right|$. By Theorem 2.2, we have
$\left|C\right|=2^{O(r)}$. So each set $U_{h}$ has cardinality at most
$r^{O(r)}\cdot r=r^{O(r)}$. Also recall that the number of sets in the MCSP
instance is
$|{\mathcal{X}}|=\sum_{j\in[r]}|{\mathcal{V}}_{j}|\cdot|L_{j}|=n^{O(r)}$, and
that the number of elements is $|U|=m=|{\mathcal{E}}|\cdot
r^{O(r)}\leq(nr)^{O(r)}=n^{O(r)}=n^{O(n^{\delta})}=2^{O(r\log r)}$. Thus
$r\geq\Omega(\log m/\log\log m)$.
The gap in the optimal congestion between the Yes-Instance and the No-Instance
is at least $r$ (see Theorem 2.2 and Lemmas 3.4, 3.10). More precisely, for
Yes-instances the congestion is at most one and for No-instances the
congestion is at least $r$. Since the approximation ratio of the alleged
algorithm is $o(\log m/\log\log m)$, it is better than $r$ for all
sufficiently large $n$ and hence the approximation algorithm can be used to
decide the satisfiability problem.
The running time of the algorithm is polynomial in the size of the MCSP
instance, i.e., is
$\mathrm{poly}(n^{O(n^{\delta})})=n^{O(n^{\delta})}=2^{O(n^{2\delta})}$. Since
$\delta>0$ is arbitrary, the theorem follows.
## 4\. Hardness of Robust k-Median on Line Metrics
We will show that the reduction from $r$-Hypergraph Label Cover to Minimum
Congestion Set Packing (MCSP) can be modified to give a $\Omega(\log
m/\log\log m)$ hardness of approximation for the Robust $k$-Median problem on
line metrics as well, where $m=|{\mathcal{S}}|$ is the number of client-sets.
For this section, it is convenient to assume that the label-sets are the
initial segments of the natural numbers, i.e.,
$L_{j}=\\{1,\ldots,\left|L_{j}\right|\\}$ and
$C=\\{1,\ldots,\left|C\right|\\}$.
Given a Hypergraph Label Cover instance $(G,\pi,r)$, we first construct a MCSP
instance $(U,{\mathcal{X}},t)$ in accordance with the procedure outlined in
Section 3.2. Next, from this MCSP instance, we construct a Robust $k$-Median
instance $(V,{\mathcal{S}},d)$ as described below.
* •
We create a location in $V$ for every set $X(v,\ell)\in{\mathcal{X}}$. To
simplify the notation, the symbol $X(v,\ell)$ will represent both a set in the
instance $(U,{\mathcal{X}},t)$, and a location in the instance
$(V,{\mathcal{S}},d)$. Thus, we have $V=\\{X(v,\ell)\in{\mathcal{X}}\\}$.
Furthermore, we create a set of clients $S(e)$ for every element $e\in U$,
which consists of all the locations whose corresponding sets in the MCSP
instance contain the element $e$. Thus, we have
${\mathcal{S}}=\\{S(e)\,:\,e\in U\\}$, where
$S(e)=\\{X(v,\ell)\in{\mathcal{X}}\,:\,e\in X(v,\ell)\\}$ for all $e\in U$.
This step is same as in Lemma 3.2.
* •
We now describe how to embed the locations in $V$ on a given line. For every
vertex $v\in{\mathcal{V}}_{j},j\in[r]$, the locations
$X(v,1),\ldots,X(v,\left|L_{j}\right|)$ are placed next to one another in
sequence, in such a way that the distance between any two consecutive
locations is exactly one. Formally, this gives
$d(X(v,\ell),X(v,\ell^{\prime}))=|\ell^{\prime}-\ell|$ for all
$\ell,\ell^{\prime}\in L_{j}$. Furthermore, we ensure that any two locations
corresponding to two different vertices in ${\mathcal{V}}$ are _not close to
each other_. To be more specific, we have the following guarantee:
$d(X(v,\ell),X(v^{\prime},\ell^{\prime}))\geq 2$ whenever $v\neq v^{\prime}$.
It is easy to verify that $d$ is a line metric.
* •
We define $k\leftarrow|{\mathcal{X}}|-t$.
Note that as $k=|{\mathcal{X}}|-t$, there is a one to one correspondence
between the solutions to the MCSP instance and the solutions to the Robust
$k$-Median instance. Specifically, a set in ${\mathcal{X}}$ is picked by a
solution to the MCSP instance iff the corresponding location is _not_ picked
in the Robust $k$-Median instance.
###### Lemma 4.1 (Yes-Instance).
Suppose that there is a labeling strategy $\sigma$ that strongly satisfies
every edge in the Hypergraph Label Cover instance $(G,\pi,r)$. Then there is a
solution to the Robust $k$-Median instance $(V,{\mathcal{S}},d)$ with
objective one.
###### Proof 4.2.
Recall the proof of Lemma 3.4. We construct a solution
${\mathcal{X}}^{\prime}\subseteq{\mathcal{X}}$, $|{\mathcal{X}}^{\prime}|=t$,
to the MCSP instance $(U,{\mathcal{X}},t)$ as follows. For every vertex
$v\in{\mathcal{V}}_{j},j\in[r]$, the solution ${\mathcal{X}}^{\prime}$
contains the set $X(v,\sigma(v))$. Now, focus on the corresponding solution
$F_{{\mathcal{X}}^{\prime}}\subseteq V$ to the Robust $k$-Median instance,
which picks a location $X$ iff $X\notin{\mathcal{X}}^{\prime}$. Hence, for
every vertex $v\in{\mathcal{V}}_{j}$, $j\in[r]$, all but one of the locations
$X(v,1),\ldots,X(v,\left|L_{j}\right|)$ are included in
$F_{{\mathcal{X}}^{\prime}}$. Since any two consecutive locations in such a
sequence are unit distance away from each other, the cost of connecting any
location in $V$ to the set $F_{{\mathcal{X}}^{\prime}}$ is either zero or one,
i.e., $d(X,F_{{\mathcal{X}}^{\prime}})\in\\{0,1\\}$ for all $X\in
V={\mathcal{X}}$.
For the rest of the proof, fix any set of clients $S(e)\in{\mathcal{S}}$,
$e\in U$. The proof of Lemma 3.4 implies that the element $e$ incurs
congestion one under ${\mathcal{X}}^{\prime}$. Hence, the element belongs to
exactly one set in ${\mathcal{X}}^{\prime}$, say $X^{*}$. Again, comparing the
solution ${\mathcal{X}}^{\prime}$ with the corresponding solution
$F_{{\mathcal{X}}^{\prime}}$, we infer that
$S(e)-F_{{\mathcal{X}}^{\prime}}=\\{X^{*}\\}$. In other words, every location
in $S(e)$, except $X^{*}$, is present in the set $F_{{\mathcal{X}}^{\prime}}$.
The clients in such locations require zero cost for getting connected to
$F_{{\mathcal{X}}^{\prime}}$. Thus, the total cost of connecting the clients
in $S(e)$ to the set $F_{{\mathcal{X}}^{\prime}}$ is at most: $\sum_{X\in
S(e)}d(X,F_{{\mathcal{X}}^{\prime}})=d(X^{*},F_{{\mathcal{X}}^{\prime}})\leq
1.$
Thus, we see that every set of clients in ${\mathcal{S}}$ requires at most
unit cost for getting connected to $F_{{\mathcal{X}}^{\prime}}$. So the
solution $F_{{\mathcal{X}}^{\prime}}$ to the Robust $k$-Median instance indeed
has objective one.
###### Lemma 4.3 (No-Instance).
Suppose that every labeling weakly satisfies at most a $2^{-\gamma r}$
fraction of the edges in the Hypergraph Label Cover instance $(G,\pi,r)$, for
some constant $\gamma$. Then every solution to the Robust $k$-Median instance
$(V,{\mathcal{S}},d)$ has objective at least $r$.
###### Proof 4.4.
Fix any solution $F\subseteq V$ to the Robust $k$-Median instance
$(V,{\mathcal{S}},d)$, and let
${\mathcal{X}}^{\prime}_{F}\subseteq{\mathcal{X}}$ denote the corresponding
solution to the MCSP instance $(U,{\mathcal{X}},t)$. Lemma 3.10 states that
there is some element $e\in U$ with congestion at least $r$ under
${\mathcal{X}}^{\prime}_{F}$. In other words, there are at least $r$ sets
$X_{1},\ldots,X_{r}\in{\mathcal{X}}^{\prime}_{F}$ that contain the element
$e$. The locations corresponding to these sets are not picked by the solution
$F$. Furthermore, the way the locations have been embedded on a line ensures
that the distance between any location and its nearest neighbor is at least
one. Hence, we have $d(X_{i},F)\geq 1$ for all $i\in[r]$. Summing over these
distances, we infer that the total cost of connecting the clients in $S(e)$ to
$F$ is at least $\sum_{i\in[r]}d(X_{i},F)\geq r$. Thus, the solution $F$ to
the Robust $k$-Median instance has objective at least $r$.
Finally, applying Lemmas 4.1, 4.3, and an argument similar to the proof of
Theorem 3.11, we get the following result.
###### Theorem 4.5.
The Robust $k$-Median problem $(V,{\mathcal{S}},d)$ is $\Omega(\log m/\log\log
m)$ hard to approximate even on line metrics, where $m=|{\mathcal{S}}|$,
unless $\textsf{NP}\subseteq\cap_{\delta>0}\textsf{DTIME}(2^{n^{\delta}})$.
## 5\. Heuristics
The Robust $k$-Median problem is a hard to approximate real-world problem and
as such heuristic solutions are interesting. In this section, we complement
our negative theoretical results with an evaluation of simple heuristics for
the Robust $k$-Median problem. In particular we look at two greedy strategies
and two variants of a local search approach. We consider a slight
generalization of the problem where clients and facilities are at separate
locations. This is more realistic and no easier than the original problem, as
one can simply place a facility at every client position to solve an instance
of the problem as defined in Definition 1.1. Due to space constraints, the
full version of this section is deferred to Appendix B.
We implemented111Code and data are available at http://resources.mpi-
inf.mpg.de/robust-k-median/code-data.7z and compared the following heuristics
to the LP relaxation (see B.1).
Greedy Upwards. Initialize all facilities as closed. Open the facility that
reduces the cost maximally. Repeat until $k$ facilities are open.
Greedy Downwards. Initialize all facilities as open. Close the facility that
increases the cost minimally. Repeat until $k$ facilities are open.
Local Search. Open $k$ random facilities. Compare all solutions that can be
obtained from the current solution by closing $\ell$ facilities and opening
$\ell$ facilities. Replace the current solution by the best solution found.
Repeat until the current solution is a local optimum. In the experiments we
use $\ell=2$.
Randomized Local Search. Same as Local Search, but instead of considering
_all_ solutions in the neighborhood, sample only a random subset. The size of
the subset is an additional parameter to the heuristic. In the experiments we
use $\ell=3$ and 200 random neighbors.
We generate three kinds of 2D-instances. In the first, _uniform_ , the clients
are uniformly distributed and all groups have the same size. The other kinds
of instances cluster the client groups according to gaussian distributions.
The intuition is that in real world instances client groups have something in
common, e.g. all come from the same city. They two kinds differ in the number
of clients per group. We have _gauss-const_ instances where all groups have
the same size and _gauss-exp_ instances where group sizes follow an
exponential distribution.
Heuristic | Uniform | Gauss-Const | Gauss-Exp
---|---|---|---
Greedy Up | 1.65 (1.49) | 5.18 (5.24) | 6.63 (5.94)
Greedy Down | 1.45 (1.42) | 2.92 (2.92) | 2.12 (2.05)
Local Search | 1.13 (1.12) | 1.63 (1.62) | 1.41 (1.39)
Randomized Local Search | 1.53 (1.48) | 2.15 (2.29) | 2.37 (2.36)
Table 1. Mean Performance as a multiple of the LP relaxation value, rounded to
three digits. In parentheses we provide the median. 1654 uniform instances,
1009 Gauss-Const instances, and 2029 Gauss-Exp instances of varying sizes were
solved. The reported performance is over the instances where the heuristics
perform worse that the LP relaxation.
Table 1 summarizes the results. The performance differences in Table 2 are
statistically significant with a very small two-sided $p$-value, according to
a Wilcoxon signed-rank test, except for the difference between Greedy Downward
and Randomized Local Search on Uniform and Gauss-Const instances. In these
cases the $p$-value is 0.66, respectively 0.08.
Since we use an LP relaxation as a comparison point, we do not know whether
the instances where the heuristics find a worse solution are actually hard for
the heuristics or whether the LP relaxation provides a much too low bound. To
investigate this we had a closer look at instances where both Greedy down and
Local Search perform badly. For three instances we solved the integer linear
program. In these instances at least it was indeed the case that the LP
relaxation yielded a bad bound. This suggests that the heuristics work even
better than the numbers in Table 2 indicate.
As expected instances where the _robust_ nature of the Robust $k$-Median
problem are not as important because groups are distributed uniformly are
easier than the more realistic instances where groups form clusters. For the
two better heuristics, Greedy Downwards and Local Search, also perform better
on instances with uneven group sizes. Here too, one can speculate that few
groups dominate the problem, and finding a solution that minimizes maximum
costs becomes easier.
The good performance of these simple heuristics indicate that although the
Robust $k$-Median problem is hard to approximate in the worst case, a
heuristic treatment can effectively find a very good approximation.
## 6\. Conclusion and Future Work
We show a logarithmic lower bound for the Robust $k$-median problem on the
uniform and line metrics, implying that there is no good approximation
algorithm for the problem. However, the empirical results suggest that real-
world instances are much easier, so it is interesting to see whether
incorporating real-world assumptions helps reducing the problem’s complexity.
For instance, if we assume that the diameter of each set $S_{i}$ is at most an
$\epsilon$ fraction of the diameter $\Delta=\max_{u,v}d(u,v)$ of the input
instance, can we obtain a constant approximation factor? This case captures
the notion of “locality” of the communities. We note that in our hardness
instances the diameter of each set $S_{i}$ is $\Delta$ for uniform metric and
at least $\Delta/2$ in the line metric, so these hard instances would not
arise if we have the locality assumption. Another interesting case is a random
instance where the sets $S_{i}$ are randomly generated by an unknown
distribution.
One can also approach this problem from the parameterized complexity angle. In
particular, can we obtain an $O(1)$ approximation algorithm in time
$g(k)\operatorname{poly}(n)$?
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## Appendix A Hypergraph Label Cover
An instance of $r$-Hypergraph Label Cover is equivalent to the $r$-Prover
system as used by Feige [8] in proving the hardness of approximation for Set
Cover. We discuss the equivalence in this section.
In the $r$-prover system, there are $r$ provers $P_{1},\ldots,P_{r}$ and a
verifier $V$. Each prover is associated with a codeword of length $r$ in such
a way that the hamming distance between any pair $P_{i},P_{j}$ is at least
${\sf ham}(P_{i},P_{j})=r/2$; this is possible if $r$ is a power of two
because we can use Hadamard code. Given an input 3-SAT formula $\phi$, the
verifier selects $r$ clauses uniformly and independently at random. Call these
clauses $C_{1},\ldots,C_{r}$. From each such clause, the verifier selects a
variable uniformly and independently at random. These variables are called
$x_{1},\ldots,x_{r}$. Prover $P_{i}$ receives a clause $C_{j}$ if the $j$th
bit of its codeword is $0$; otherwise, it receives variable $x_{j}$. The
property of Hadamard code guarantees that each prover would receive $r/2$
clauses and $r/2$ variables.
Then each prover $P_{i}$ is expected to give an assignment to all involved
variables it receives and sends this assignment to the verifier. The verifier
then looks at the answers from $r$ provers and has two types of acceptance
predicates.
* •
(Weak acceptance) At least one pair of answers is consistent.
* •
(Strong acceptance) All pairs of answers are consistent.
Applying parallel repetition theorem [15], Feige argues the following.
###### Theorem A.1.
([8, Lemma 2.3.1]) If $\Phi$ is a satisfiable 3-SAT(5) formula, then there is
provers’ strategy that always causes the verifier to accept. Otherwise, the
verifier weakly accepts with probability at most $r^{2}2^{-\gamma r}$ for some
universal constant $\gamma>0$.
Now we show how Theorem 2.2 follows by constructing the instance of Hypergraph
Label Cover $(V,E)$ based on the $r$-prover system. For each prover $j$, we
create a set $V_{j}$ consisting of vertices $v$ that correspond to possible
query sent to prover $j$, so we have $|V_{j}|=(5n/3)^{r/2}n^{r/2}$. For each
possible random string $x$, we have an edge $h_{x}$ that contains $r$
vertices, corresponding to queries sent to the provers. It can be checked that
the total number of possible random strings is $(5n)^{r}$, and the degree of
each vertex is $3^{r/2}5^{r/2}=15^{r/2}$; notice that this is equal to
$r|E|/|V|$. A prover strategy corresponds to the label of vertices, and the
acceptance probability is exactly the fraction of satisfied edges. Moreover,
for each possible query, the number of possible answers is at most $7^{r}$
(for each clause, there are $7$ ways to satisfy it). This implies that
$|L_{j}|\leq 7^{r}$.
## Appendix B Heuristics
The Robust $k$-Median problem is a real-world problem and as such needs to be
solved as well as possible despite its hardness of approximation. In this
section, we complement our negative theoretical results with an experimental
evaluation of different simple heuristics for the Robust $k$-Median problem.
In particular we look at two variants of a greedy strategy and two variants of
a local search approach. We consider a slight generalization of the problem
where clients and facilities are at separate locations. This is more realistic
and no easier than the original problem, as one can simply place a facility at
every client position to solve an instance of the problem as defined in
Definition 1.1.
This is by no means an exhaustive exploration of the possible solution space.
However, the results we obtain indicate that a heuristic treatment of the
Robust $k$-Median problem can yield surprisingly good solutions, even if the
heuristics are very naive.
For our experiments we consider instances in the plane, as these are closest
to the real-world motivation for the problem. We wanted to check how the
structure of the instance influences the performance of the heuristics. We
suspected that instances where the clients are distributed uniformly are easy,
as intuitively a solution that is good for one group of clients is good for
all groups.
The robust version of the k-median problem is considered because often the
exact set of clients is not known before choosing facility locations and one
wants to perform well even if the worst set of possible clients turns out to
be realized. It is reasonable to assume that every group of clients has
something in common, for example that they come from a similar region, like a
city. Therefore more realistic instances for the Robust $k$-Median problem
have the groups form clusters in space. We also generate such instances for
testing our heuristics.
### B.1. Methods
Since solving Robust $k$-Median instances to optimality is infeasible for the
instances we consider222We attempted solving three instances optimally, see
Figure 2, but gave up on the third after nearly half a year of CPU time was
consumed., we compare the performance of the various heuristics to the value
of a LP-relaxation. We have a variable $x_{j}$ for each possible median
location and variables $y_{ij}$ that indicate whether client $i$ is served by
facility $j$. The LP is then as follows.
$\displaystyle\min\quad$ $\displaystyle T$ $\displaystyle s.t.\quad$
$\displaystyle y_{ij}-x_{j}\leq 0$ $\displaystyle\forall i,j$
$\displaystyle\sum_{j}y_{ij}\geq 1$ $\displaystyle\forall i$
$\displaystyle\sum_{i\in g}d(i,j)\cdot y_{ij}\leq T$
$\displaystyle\forall\text{ groups of clients }g$
$\displaystyle\sum_{j}x_{j}\leq k\quad\text{and}\quad 0\leq x_{j}\leq 1$
$\displaystyle\forall j$ $\displaystyle 0\leq y_{ij}\leq 1$
$\displaystyle\forall i,j$
To solve the LP we use the Gurobi solver [9], version 5.5.0, on a 64-bit Linux
system.
Note that the assignment of the $y_{ij}$ variables is immediately clear from
the assignment of the $x_{j}$. For location $i$, let $j_{1}$, $j_{2}$,
…$j_{n}$ be the locations ordered by increasing distance. Then
$y_{ij_{\ell}}=\min(x_{j_{\ell}},1-(y_{ij_{1}}+\ldots+y_{ij_{\ell-1}}))$. The
constraint $y_{ij_{\ell}}\leq\min(\,,\,)$ is already expressed by the first
two constraints. It could however be put into the objective via the big
$M$-method. Consider a minimization problem $\min T$ subject to $x=\min(b,c)$.
Let $M$ be large integer and consider $\min T+Mt$ subject to $x\leq b$, $x\leq
c$, $t\leq b-x$, and $t\leq c-x$. Observe that $t=\min(b,c)-x$ in an optimal
solution. One needs to choose $M$ big enough so that $t$ must be zero in an
optimal vertex solution. It is however unclear whether this will speed up the
solution. We have not tried this method.
We implemented and compared the following heuristics:
Greedy Upwards. Initialize all facilities as closed. Open the facility that
reduces the cost maximally. Repeat until $k$ facilities are open.
Greedy Downwards. Initialize all facilities as open. Close the facility that
increases the cost minimally. Repeat until $k$ facilities are open.
Local Search. Open $k$ random facilities. Compare all solutions that can be
obtained from the current solution by closing $\ell$ facilities and opening
$\ell$ facilities. Replace the current solution by the best solution found.
Repeat until the current solution is a local optimum. In the experiments we
use $\ell=2$.
Randomized Local Search. Same as Local Search, but instead of considering
_all_ solutions in the neighborhood, sample only a random subset. The size of
the subset is an additional parameter to the heuristic. In the experiments we
use $\ell=3$ and 200 random neighbors.
Note that taking the solution of one of the greedy algorithms as starting
point for a local search is an obvious improvement, but this would prevent us
from comparing the local search algorithm with the greedy heuristic.
The local search heuristic is closely related to Lloyd’s algorithm for the
k-means problem. In Lloyd’s algorithm, a random set of centers is chosen and
iteratively updated by moving the centers to the centroids of the clients that
fall in their voronoi cell. This improves the total distance from the centers
to all clients in every iteration.
In our setting, we want to reduce the cost of the group of clients that
currently incurs the maximal cost. This can be done by moving a facility
closer to this group of clients, that is, closing one facility and opening
another that reduces the objective function. The local search algorithm, by
closing and opening more than one facility at a time, does this at least as
well.
We create instances in the plane and use the euclidean distance. We create two
types of instances. In the first type the clients and facilities are uniformly
distributed in a $100\times 100$ square. We call these instances the _uniform_
instances. In these instances all groups of clients contain the same number of
clients. The $k$ we use for the experiments is 7.
The second kind uses random gaussian distributions to sample client positions.
To generate the gaussian distributions we sample a matrix $M$ with
$v_{1},v_{2}$ on the diagonal, where the two values are chosen uniformly at
random from $[0,50]$, the matrix is then rotated by a uniformly random angle.
The result is the covariance matrix of the gaussian distribution. The mean is
a random point in a $100\times 100$ square. These instances we call _gauss_.
We generate two subgroups of instances, in the first subgroup, _gauss-const_ ,
all groups of clients have the same number of clients, in the second subgroup,
_gauss-exp_ , the number of clients in a group is sampled from an exponential
distribution. Figure 1 shows examples for the different kind of instances we
generate.
(a) Uniform
(b) Gauss-Const
(c) Gauss-Exp
Figure 1. Examples for the kind of instances we generate. Circles are clients,
squares are facilities, colors indicate group membership.
As we didn’t put much effort into optimizing our heuristics for speed (for
example we don’t use spatial search structures to find nearest neighbors), we
don’t report execution time and focus solely on solution quality. Nevertheless
it is clear that the greedy strategies are much simpler to implement and much
faster than the local search heuristics.
We report average performance on instances where the solution is worse than
the LP value, as small, easy instances otherwise skew the results. To conclude
relative performance advantages between heuristics we use a Wilcoxon signed-
rank test as implemented in SciPy 0.12.0.
All computer code we wrote to run the experiments and analyze the results, as
well as the instances we solved, is available online at http://resources.mpi-
inf.mpg.de/robust-k-median/code-data.7z.
### B.2. Results
Heuristic | Uniform | Gauss-Const | Gauss-Exp
---|---|---|---
Greedy Up | 1.65 (1.49) | 5.18 (5.24) | 6.63 (5.94)
Greedy Down | 1.45 (1.42) | 2.92 (2.92) | 2.12 (2.05)
Local Search | 1.13 (1.12) | 1.63 (1.62) | 1.41 (1.39)
Randomized Local Search | 1.53 (1.48) | 2.15 (2.29) | 2.37 (2.36)
Table 2. Mean Performance as a multiple of the LP relaxation value, rounded to
three digits. In parentheses we provide the median. 1654 uniform instances,
1009 Gauss-Const instances, and 2029 Gauss-Exp instances of varying sizes were
solved. The reported performance is over the instances where the heuristics
perform worse that the LP relaxation.
Clients | Facilities
---|---
| 10 | 110 | 210 | 310 | 410
| Greedy | Search | GD | LS | GD | LS | GD | LS | GD | LS
10 | 1.00 | 1.00 | 1.12 | 1.00 | 1.31 | 1.01 | 1.39 | 1.02 | 1.40 | 1.01
160 | 1.01 | 1.01 | 1.6 | 1.17 | 1.63 | 1.17 | 1.68 | 1.15 | 1.63 | 1.15
310 | 1.01 | 1.01 | 1.64 | 1.21 | 1.69 | 1.19 | 1.70 | 1.19 | 1.75 | 1.18
460 | 1.01 | 1.01 | 1.68 | 1.22 | 1.73 | 1.21 | 1.71 | 1.21 | 1.73 | 1.21
110 | 1.00 | 1.00 | 1.17 | 1.01 | 1.22 | 1.01 | 1.25 | 1.01 | 1.24 | 1.01
1760 | 1.0 | 1.0 | 1.28 | 1.06 | 1.33 | 1.06 | 1.34 | 1.06 | 1.34 | 1.06
3410 | 1.0 | 1.0 | 1.3 | 1.07 | 1.33 | 1.07 | | | |
(a) Uniform
Clients | Facilities
---|---
| 10 | 110 | 210 | 310 | 410
| Greedy | Search | GD | LS | GD | LS | GD | LS | GD | LS
10 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0
160 | 1.0 | 1.0 | 2.74 | 1.64 | 3.05 | 1.6 | 3.33 | 1.62 | 3.33 | 1.57
310 | 1.0 | 1.0 | 2.76 | 1.70 | 3.07 | 1.66 | 3.32 | 1.64 | |
110 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0∗ | 1.0∗ | | | |
3410 | 1.01 | 1.0 | 2.74 | 1.65 | 3.02∗ | 1.63∗ | | | |
(b) Gauss-Const
Clients | Facilities
---|---
| 10 | 110 | 210 | 310 | 410
| Greedy | Search | GD | LS | GD | LS | GD | LS | GD | LS
10 | | | | | 1.0∗ | 1.0∗ | 1.0 | 1.0 | 1.0 | 1.0
110 | 1.0 | 1.0 | 1.34 | 1.16 | 1.66 | 1.28 | 1.65 | 1.26 | 1.91 | 1.34
210 | 1.0 | 1.0 | 1.9 | 1.41 | 2.14 | 1.45 | 2.31 | 1.46 | 2.46 | 1.49
310 | 1.0 | 1.0 | 2.23 | 1.48 | 2.6 | 1.48 | 2.69 | 1.51 | 2.78 | 1.50
110 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0 | 1.0 | 1.01 | 1.0
1210 | 1.0 | 1.0 | 1.38 | 1.21 | 1.56 | 1.23 | 1.73 | 1.29 | 1.77 | 1.29
2310 | 1.0 | 1.0 | 1.94 | 1.38 | 2.09 | 1.44 | 2.48 | 1.41 | 2.29 | 1.44
3410 | 1.0 | 1.0 | 2.17 | 1.51 | 2.48 | 1.48 | 2.8 | 1.55 | |
(c) Gauss-Exp
Table 3. Performance depending on instance size for the Greedy Downwards and
Local Search heuristics. All values are averages over 50 instances, except for
those marked by ∗. For Gauss-Exp instances the number of clients is the mean
of the exponential distribution times the number of groups. Values above the
horizontal line come from instances with 10 clients per group, below the line
instances have 110 clients per group.
Table 2 summarizes the results of the experiments, Table 3 shows the
performance for the different instance sizes for the Greedy Upwards and the
Local Search heuristic. The performance differences in Table 2 are
statistically significant with a very small two-sided $p$-value, except for
the difference between Greedy Downward and Randomized Local Search on Uniform
and Gauss-Const instances. In these cases the $p$-value is 0.66, respectively
0.08.
Since we use an LP relaxation as a comparison point, we do not know whether
the instances where the heuristics find a worse solution are actually hard for
the heuristics or whether the LP relaxation provides a much too low bound. To
investigate this we had a closer look at instances where both Greedy down and
Local Search perform badly. For three instances we attempted to solve the
integer linear program and succeeded for two of them. In Figure 2 we see
different solutions. For these instances at least it was indeed the case that
the LP relaxation yielded a bad bound. This suggests that the heuristics work
even better than the numbers in Table 2 indicate.
(a) Uniform: LP value 2806.4, Greedy value 5982.39, Local Search value
3426.43, OPT 3230.19.
(b) Gauss-Const: LP value 1360.26, Greedy value 8307.48, Local Search value
2541.21, OPT 2505.26
(c) Gauss-Exp: LP value 2362.06, Greedy value 10624.4, Local Search value
4354.54, 4192.31 $\leq$ OPT $\leq$ 4354.54
Figure 2. Solutions of the different algorithms on particularly hard
instances. From left to right, the LP solution, the Greedy downwards solution,
the Local Search solution and the ILP solution. Darkness of facilities
indicates "how open" they are in the LP relaxation. In 2c the ILP solver was
still running at the time of submission, after having consumed 177 days of CPU
time and 46GB of memory.
### B.3. Conclusion
Note that all heuristics perform very well on the instances we tried. In
accordance with our theoretical results, increasing the number of groups makes
the instances harder, more so that increasing the number of facilities or the
number of clients.
As expected instances where the _robust_ nature of the Robust $k$-Median
problem are not as important because groups are distributed uniformly are
easier than the more realistic instances where groups form clusters. For the
two better heuristics, Greedy Downwards and Local Search, also perform better
on instances with uneven group sizes. Here too, one can speculate that few
groups dominate the problem, and finding a solution that minimizes maximum
costs becomes easier.
The good performance of these simple heuristics indicate that although the
Robust $k$-Median problem is hard to approximate in the worst case, a
heuristic treatment can effectively find a very good approximation.
|
arxiv-papers
| 2013-09-18T10:11:31 |
2024-09-04T02:49:51.078660
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Sayan Bhattacharya, Parinya Chalermsook, Kurt Mehlhorn, Adrian Neumann",
"submitter": "Adrian Neumann",
"url": "https://arxiv.org/abs/1309.4602"
}
|
1309.4612
|
The Innermost Regions of Relativistic Jets and Their Magnetic Fields
11institutetext: Departamento de Astronomía y Astrofísica, Universidad de
Valencia, SPAIN
# Numerical study of broadband spectra caused by internal shocks in magnetized
relativistic jets of blazars.
Jesús M. Rueda-Becerril 11 [email protected] Petar Mimica 11 Miguel A. Aloy
11 Carmen Aloy 11
###### Abstract
The internal-shocks scenario in relativistic jets has been used to explain the
variability of blazars’ outflow emission. Recent simulations have shown that
the magnetic field alters the dynamics of these shocks producing a whole zoo
of spectral energy density patterns. However, the role played by magnetization
in such high-energy emission is still not entirely understood. With the aid of
_Fermi_ ’s second LAT AGN catalog, a comparison with observations in the
$\gamma$-ray band was performed, in order to identify the effects of the
magnetic field.
## 1 Introduction
Relativistic outflows have been observed extensively in blazars, a class of
radio-loud active galactic nuclei (AGNs) whose jets are pointing very close to
the line of sight towards the observer Urry:1995aa , and known for showing the
most rapid variability of all AGNs. Their remarkable characteristic flares in
the X-ray frequency range usually have a duration of the order of one day.
Often the internal-shocks (IS) scenario Rees:1994ca is invoked to explain
this variability Spada:2001do ; Mimica:2004ay ; Mimica:2005aa . The IS
scenario is an idealized model of a variable jet where an intermittently
working central engine ejects shells of magnetized plasma which collide due to
their velocity differences. As a consequence of the collision internal shocks
are formed, particles are accelerated at the shock fronts and the non-thermal,
highly variable radiation is produced.
Our long-term project is the study of the influence of magnetic fields on the
radiation from IS using numerical simulations. In Mimica:2012aa we studied a
large number of shell collisions with different magnetization levels. In the
present work we focus on a limited number of shell magnetization levels, but
vary other parameters such as the jet viewing angle, bulk Lorentz factor of
the shells, and their relative Lorentz factor. The data obtained from these
simulations is used to categorize the specific effects that variations of each
parameter have on average spectra. These synthetic observations are then
compared with the second LAT AGN catalog (2LAC) of blazars observed by Fermi
Ackermann:2011apj .
Our numerical setup is describen in Sec. 2. The results are shown in Sec. 3.
Finally we discuss briefly our results in Sec. 4.
## 2 Numerical Setup
We use a modified version of the _SPEV_ codeMimica:2009aa ; Mimica:2012aa to
compute the non-thermal emission from the IS. We do not consider the full
hydrodynamic interaction of colliding magnetized shells (see. e.g.
Mimica:2007aa for a detailed study). Instead, we simplify the shell
interaction as a one-dimensional Riemann problem and focus our resources on a
more detailed treatment of the non-thermal radiation. Our method consists of
three phases:
1. 1.
_Solution of the Riemann problem._ Making use of an exact RMHD Riemann solver
Romero:2005zr we determine the properties of the internal shock waves. We
follow the procedure described in Mimica:2010aa to set-up the shells and to
extract the information needed for the steps 2 and 3.
2. 2.
_Non-thermal particles transport and evolution._ The particles are injected
behind the shock fronts following the prescription of Bottcher:2010gn ;
Joshi:2011bp ; Mimica:2012aa . We assume that a fraction of the thermal
electrons are accelerated to high energies, and that their energy density is a
fraction of the internal energy density of the shocked fluid. We assume a
cylindrical shell geometry and perform all the calculations in the rest frame
of the shocked fluid. In this frame the shocks are propagating away from the
initial discontinuity, injecting and leaving non-thermal particles behind. We
evolve the energy distribution of non-thermal electrons taking into account
synchrotron and inverse-Compton (IC) losses. See Mimica:2012aa for more
details.
3. 3.
_Radiative transfer._ The total emissivity at each point is assumed to be a
combination of the following emission processes: (1) synchrotron radiation,
(2) IC upscattering of an external radiation field (EIC) and (3) synchrotron
self-Compton (SSC). The details of how they are calculated are given in
Mimica:2012aa . The radiative transfer equation is solved taking into account
the relativistic effects and time delays.
We compute light curves and average spectral energy distribution (SED) for
each shell collision. In this work we focus our attention on how parameter
variations affect the SEDs.
## 3 Results
As mentioned in Sec. 1, the aim of this work is to cover a wider range in the
parameter space than was done in Mimica:2012aa . We group our models according
to the initial shell magnetization, $\sigma:=B^{2}/4\pi\rho\Gamma^{2}c^{2}$.
We denote by letters S, M and W the following families of models:
* W:
weakly magnetized, $\sigma_{L}=10^{-6},\sigma_{R}=10^{-6}$,
* M:
moderately magnetized, $\sigma_{L}=10^{-2},\sigma_{R}=10^{-2}$, and
* S:
strongly magnetized, $\sigma_{L}=1,\sigma_{R}=10^{-1}$.
Hereafter the subscripts $L$ and $R$ will denote left (faster) and right
(slower) shells, respectively. As parameters to vary we considered both
intrinsic and extrinsic ones. Among the intrinsic parameters we choose the
Lorentz factor of the slow shell, $\Gamma_{R}$, and the relative Lorentz
factor $\Delta g:=\Gamma_{L}/\Gamma_{R}-1$, where $\Gamma_{L}$ is the Lorentz
factor of the fast shell. The parameter space covered is shown in Table 1.
Parameter | value
---|---
$\sigma_{L}$ | $10^{-6},\ 10^{-2},\ 1$
$\sigma_{R}$ | $10^{-6},\ 10^{-2},\ 10^{-1}$
$\Gamma_{R}$ | $10,12,17,20,22,25$
$\Delta g$ | $0.5,0.7,1.0,1.5,2.0$
$\theta$ | $5$
Table 1: Parameters of the models. $\Gamma_{R}$ is the Lorentz factor of the
slow shell, $\sigma_{L}$ and $\sigma_{R}$ are the fast and slow shell
magnetizations and $\theta$ is viewing angle of the observer.
For clarity, when we refer to a particular model we label it by appending the
values of each of these parameters to the model letter. For instance,
S-G10-D1.0-T5 is the strongly magnetized model with $\Gamma_{R}=10$ (G10),
$\Delta g=1.0$ (D1.0) and $\theta=5^{\circ}$ (T5). If we refer to a subset of
models with one or two parameters fixed we use an abbreviated notation, where
we omit the varying parameters from the label. We compute the spectra for a
typical source located at $z=0.5$.
In the rest of this section we will present some of the final SEDs resulting
from our simulations. A larger collection is shown in (Rueda:2013aa, ). The
SEDs of each model has been averaged over the time interval $0-10^{6}$ s.
### 3.1 Weakly-magnetized models
The SEDs computed for the models W-G10-T5 (varying $\Delta g$) are shown in
the left panel of Fig. 1. The spectra show that with increasing $\Delta g$ the
IC component also increases, up to three orders of magnitude. In order to see
the effects on each emission process, the synchrotron, SSC and EC components
for $\Delta g=0.5,2.0$ are shown as dashed, dot-dashed and dot-dot-dashed
lines, respectively. As we can see, while the three components of the spectrum
(synchrotron, SSC and EC) are around the same order of magnitude for $\Delta
g=0.5$, for $\Delta g=2.0$ the SSC is almost two orders of magnitude more
luminous than the other two. The inset shows the $\gamma$-ray spectral slope
of each model as a function of its $\gamma$-ray flux (see Sec. 3.4).
### 3.2 Moderately-magnetized models
The SEDs of the family of models M-D1.0-T5 are presented in the right panel of
Fig. 1. Analogous to the left panel of Fig. 1 the synchrotron, SSC and EC
components are shown as dashed, dot-dashed and dot-dot-dashed lines,
respectively, for $\Gamma_{R}=10,17,25$. The synchrotron component for
$\Gamma_{R}=10$ is $\simeq 20$ times brighter than the SSC one, in contrast to
the EC which is $100$ times dimmer. For $\Gamma_{R}=25$ the EC is of the same
order of magnitude of SSC and synchrotron. The latter two decrease one order
of magnitude between M-G10-D1.0-T5 and M-G25-D1.0-T5, while the EC grows by
almost one order of magnitude. This is a consequence of the fact that the
number of electrons and the comoving magnetic field strength decrease with the
increasing $\Gamma_{R}$ Mimica:2012aa , which means that there are less
synchrotron photons and less electrons which can scatter them in the SSC
process. On the other hand, the radiation field density of the seed photons
for the EC is independent of $\Gamma_{R}$, which, in combination with the
Doppler boost causes the increase in the EC luminosity. We also see that at
$\simeq 10^{23}$ Hz there is a point where all the EC spectra coincide. This
is due to the Klein-Nishina cutoff, which we model as a sharp cutoff. The
inset shows that the there is no significant change in the flux of
$\gamma$-ray photons, although there was for the spectral index, heading
towards lower values for increasing $\Gamma_{R}$.
Figure 1: Left panel: Averaged spectra resulting from the collision of weakly
magnetized shells ($\sigma_{L}=\sigma_{R}=10^{-6}$). For the models $\Delta
g=0.5,2.0$, the synchrotron, SSC and EIC contributions (dashed, dot-dashed and
dot-dot-dashed lines, respectively) are shown. The inset shows the photon
spectral slope $\Gamma_{ph}$ as a function of the photon flux $F_{ph}$ in the
$\gamma$-ray band (see Sec. 3.4). Colors of the points correspond to the line
colors in the main plot. Right panel: Same as the left panel, but for
moderately magnetized shells ($\sigma_{L}=\sigma_{R}=10^{-2}$) and varying
$\Gamma_{R}$. For models with $\Gamma_{R}=10,17,25$ we plot the synchrotron,
SSC and EIC contributions.
### 3.3 Strongly-magnetized models
The third model family consists of the strongly magnetized models where
$\sigma_{L}=1$ and $\sigma_{R}=0.1$. The SEDs of the series of models
S-D1.0-T5 appear in Fig. 2. As we can see, for $\Gamma_{R}=10$ the synchrotron
component is $\simeq 100$ times brighter than the IC. For $\Gamma_{R}=25$ this
difference decreases to one order of magnitude. The EIC component rises with
rising $\Gamma_{R}$, to the point in which it begins to be comparable to the
synchrotron component. These effects are similar to the family M-D1.0-T5,
described in Sec. 3.2. The spectra converge due to our treatment of the Klein-
Nishina cutoff. In the inset we can see that the flux of $\gamma$-ray photons
does not change appreciably in this family of models.
Figure 2: Same as right panel in Fig. 1, but for strongly magnetized shells
($\sigma_{L}=1,\sigma_{R}=$ 0.1).
### 3.4 $\gamma$-rays spectral slope
A linear mean-squares algorithm is used to deduce the $\gamma$-ray spectral
slope $\Gamma_{ph}$. Due to the fact that we are not modeling the Klein-
Nishina part of the spectrum, we only performed the calculations of
$\Gamma_{ph}$ for those models that do not show a large drop-off in the photon
flux. In Fig. 3 we show $\Gamma_{ph}$ as a function of the photon flux for
energies $>0.2\,$GeV, where $F_{ph}$ is the photon flux for photon energies
$>0.1$ GeV Abdo:2009cb . We compare our results with sources found in 2LAC
catalogue Ackermann:2011apj (restricting the comparsion to sources with
$0.4\leq z\leq 0.6$). In Fig. 3 we see that weakly and moderately magnetized
models overlap with the observations, with more weakly than moderately
magnetized models falling within the observed part of the parameter space.
Preliminary results of models where the viewing angle, $\theta$, is changed;
i. e. SMW-G10-D1.0, appear also in Fig. 3.
Figure 3: Spectral slope $\Gamma_{ph}$ as a function of the photon flux for
energies $>100$ MeV. $\Gamma_{ph}$ is computed for the photon energies $>200$
MeV Abdo:2009cb . The symbols joined by lines represent our numerical models,
while cyan and magenta triangles represent BL Lacs and FSRQs at redshift
$z\simeq 0.5$ from 2LAC Ackermann:2011apj . In this figure we also show the
preliminary three families of models where we vary the opening angle (filled
circles).
## 4 Conclusions
In this paper we report on the progress of the study of the influence of the
jet magnetization on blazar flares. We vary two parameters of our models: the
relative Lorentz factor $\Delta g$ and the bulk Lorentz factor $\Gamma_{R}$.
When $\Delta g$ is increased we get a more luminous maximum of the inverse
Compton component, which is dominated by the SSC. If $\Gamma_{R}$ is increased
we find that the EIC begins to dominate over SSC, as well as becoming
comparable to the synchrotron component. In the case of strongly magnetized
shells, if $\Gamma_{R}\sim 50$ both synchrotron and EC components are of the
same order of magnitude (see Rueda:2013aa ). Among all the models studied
here, the weakly magnetized are the ones that best fit _Fermi_ observations
Ackermann:2011apj . However, the tendencies of certain models with higher
magnetization appear to also be consistent with the observations.
## Acknowledgments
JMRB acknowledges the support from the Grisolia fellowship GRISOLIA/2011/041.
PM, MAA and CA acknowledge the support from the ERC grant CAMAP-259276 and the
grants AYA2010-21097-C03-01 and PROMETEO-2009-103.
## References
* (1) C.M. Urry, P. Padovani, PASP 107, 803 (1995)
* (2) M.J. Rees, P. Meszaros, ApJL 430, L93 (1994)
* (3) M. Spada, G. Ghisellini, D. Lazzati, A. Celotti, MNRAS 325, 1559 (2001)
* (4) P. Mimica, M.A. Aloy, E. Müller, W. Brinkmann, A&A 418, 947 (2004)
* (5) P. Mimica, M.A. Aloy, E. Müller, W. Brinkmann, A&A 441, 103 (2005)
* (6) P. Mimica, M.A. Aloy, MNRAS 421, 2635 (2012)
* (7) M. Ackermann et al., ApJ 743, 171 (2011)
* (8) P. Mimica, M.A. Aloy, I. Agudo, J.M. Marti, J.L. Gómez, J.A. Miralles, ApJ 696, 1142 (2009)
* (9) P. Mimica, M.A. Aloy, E. Müller, A&A 466, 93 (2007)
* (10) R. Romero, J. Marti, J.A. Pons, J.M. Ibáñez, J.A. Miralles, JFM 544, 323 (2005)
* (11) P. Mimica, M.A. Aloy, MNRAS 401, 525 (2010)
* (12) M. Böttcher, C. Dermer, ApJ 711, 445 (2010)
* (13) M. Joshi, M. Böttcher, ApJ 727, 21 (2011)
* (14) J. Rueda-Becerril, P. Mimica, M.A. Aloy, MNRAS, in preparation (2013)
* (15) A.A. Abdo, M. Ackermann, M. Ajello, W.B. Atwood, M. Axelsson, L. Baldini, J. Ballet, B. et al., ApJ 700, 597 (2009), 0902.1559
|
arxiv-papers
| 2013-09-18T11:09:03 |
2024-09-04T02:49:51.090005
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Jesus M. Rueda-Becerril, Petar Mimica, Miguel A. Aloy, Carmen Aloy",
"submitter": "Jes\\'us Misr\\'ayim Rueda-Becerril",
"url": "https://arxiv.org/abs/1309.4612"
}
|
1309.4630
|
# Radio constraints on Galactic WIMP dark matter
Roberto A. Lineros Instituto de Física Corpuscular (CSIC-Universitat de
València), Apdo. 22085, E-46071 Valencia, Spain. [email protected]
###### Abstract
Synchrotron emission from electron cosmic ray populations can be used to study
both cosmic rays physics and WIMP dark matter imprints on radio skymaps. We
used available radio data - from MHz to GHz - to analyze the contribution from
galactic WIMP annihilations and impose constraints on WIMP observables:
annihilation cross section, channel and mass. Depending on the annihilation
channel we obtain as competitive bounds as those obtained in FERMI-LAT gamma
ray analysis of dwarf satellite galaxies.
###### keywords:
dark matter theory , cosmic ray theory , absorption and radiation processes
††journal: NIM A RICAP-2013
## 1 Introduction
During the last decades improvements on cosmological observations have
confirmed that around 27% of the matter content of observable Universe is
dominated by non-barionic Dark Matter (DM) [1]. The most popular DM candidates
are in the form of Weakly Interactive Massive Particles (WIMPs) due to their
diverse and complementary ways of detection. Moreover, WIMPs are common to
many particle physics models beyond the Standard Model (SM) e.g.
Supersymmetry, extra-dimensions, and extended Higgs sector models. All these
candidates share similar characteristics: non-colored, electrically neutral
and stable on cosmological scales. WIMPs participate actively in the early
universe until they freeze out from the thermal bath due to the expansion of
the universe. This mechanism provide a natural way to reproduce the observed
relic abundance [1],
$\Omega_{\rm DM}h^{2}=0.1196\pm 0.0031\,,$ (1)
and it requires cross section of the order of electroweak interactions. This
cross section is usually referred as the thermal WIMP cross section:
$\langle\sigma v\rangle=3\times 10^{-26}\,{\rm cm}^{3}{\rm s}^{-1}.$
One of the expected signatures are the annihilation products from WIMPs around
celestial objects like the Milky Way and its dwarf satellite galaxies. The SM
particles coming from WIMP annihilations can produce gamma rays, cosmic rays,
and neutrinos. Current and projected observatories can explore these new
signals thanks to their high sensitivity and resolution. For example, this is
the case of the FERMI-LAT gamma ray observatory. On the other hand, radio
surveys have been continuously observing and characterizing the galactic
activity. One of the most known all–sky survey is the HASLAM map [2] at the
frequency of 408 MHz and it is widely used in astrophysics.
This work is based on Ref. [3] and describes its main results.
## 2 Constraints from radio data
We consider the synchrotron emission at radio frequencies, i.e. from MHz to
GHz, from electrons and positrons as result of WIMP annihilations of the DM
halo of the Milky Way. In order to cover the whole possibilities, we take into
account the uncertainties regarding the production of synchrotron emission:
* (i)
Galactic magnetic field (GMF) ranging between 1–10 $\mu$G.
* (ii)
Physically motivated GMF spatial distributions.
* (iii)
Cosmic ray propagation models.
* (iv)
DM distributions: isothermal and NFW profile [4].
* (v)
DM annihilation into: muons, taus, $W$ bosons and $b$ quarks.
* (vi)
DM annihilation cross section similar to the thermal value.
For most of these scenarios, the radio emission coming from WIMPs is as bright
as the available radio surveys. Thus, using the radio survey we set
constraints for some of the DM properties e.g. the annihilation cross section.
The bounds are then obtained by requiring that the DM signal would not exceed
the observed radio emission by 3 standard deviations, i.e.
$T_{\rm DM}(\nu)\leq T_{\rm obs}(\nu)+3\,\sigma_{\rm survey}(\nu)\,.$ (2)
The radio surveys and the DM emision skymap where previously divided into
patches of $\sim 10^{\circ}\times 10^{\circ}$ and then compared. We did not
include any astrophysical radio sources meaning that our bounds are
conservative.
Figure 1: Upper bounds on the annihilation cross section versus DM mass for
annihilations into $b$. In dashed line, we show bounds obtained by FERMI-LAT
from the analysis on dwarf satellite galaxies [5]. Radio constraints are less
competitive with respect to hadronic channels. Further details in Ref. [3].
Figure 2: Upper bounds on the annihilation cross section versus DM mass for
annihilations into $W$. Blueish dashed line are the bounds from FERMI-LAT
analysis, also we show the bounds from antiproton/proton analysis [6]. Further
details in Ref. [3]. Figure 3: Upper bounds on the annihilation cross section
versus DM mass for annihilations into muons. Radio constraints extracted from
leptonic channel are in better situation than constraints obtained by FERMI-
LAT. Further details in Ref. [3].
This procedure mainly constrains the annihilation cross section for a specific
annihilation channel. In Figs. 1, 2 and 3, we present some representative
results of the analysis carried in Ref. [3]. All of these show upper bounds on
the annihilation cross section for different annihilation channels versus the
WIMP mass. In Fig. 1, we analyze the case of annihilation into $b$ quarks. We
note that our bounds are less restrictive than those obtained in the gamma ray
analysis of dwarf satellite galaxies [5]. This is mainly due to annihilations
into quarks produce much more gamma rays because of large $\pi^{0}$ production
and subsequent decay. Similar situation happens in Fig. 2 when we compare our
results with previous works on constraints from antiproton/proton observation
[6].
The opposite occurs when we compare the muon channel, Fig. 3. In this case,
radio bounds are more competitive than those from [5] because muonic
annihilation channel (and in general any leptonic channel) produce larger
amount of electrons and positrons.
For the $b$ and muon channels, the bounds reach the thermal cross section
value at low WIMP masses. This method constrains particle physics models with
predicted annihilation channels and cross section higher than our bounds.
Our analysis shows the complementarity between radio and gamma rays
observations for the study of galactic DM. In principle, a combined analysis
in this direction would improve DM indirect searches.
## 3 Conclusions
Synchrotron emission from galactic WIMP annihilations presents an interesting
and alternative observable to study the WIMP properties. This method is
specially efficient for constraining annihilation into leptons. Our bounds are
quite competitive with respect to similar analysis and reach the value of WIMP
thermal cross section for WIMP masses lower than 15 GeV.
## Acknowledgements
We acknowledge N. Fornengo, M. Regis and M. Taoso. This work was supported by
the Spanish MINECO under grants FPA2011-22975 and MULTIDARK CSD2009-00064
(Consolider-Ingenio 2010 Programme). Also, it was supported by
Prometeo/2009/091 (Generalitat Valenciana), and by the EU ITN UNILHC PITN-
GA-2009-237920.
## References
* [1] Planck Collaboration, Planck 2013 results. XVI. Cosmological parameters arXiv:1303.5076.
* [2] C. G. T. Haslam, C. J. Salter, H. Stoffel, W. E. Wilson, A 408 MHz all-sky continuum survey. II - The atlas of contour maps, A&A Supplement 47 (1982) 1.
* [3] N. Fornengo, R. A. Lineros, M. Regis, M. Taoso, Galactic synchrotron emission from WIMPs at radio frequencies, JCAP1 (2012) 5. arXiv:1110.4337, doi:10.1088/1475-7516/2012/01/005.
* [4] J. F. Navarro, C. S. Frenk, S. D. M. White, The Structure of Cold Dark Matter Halos, Astrophys. J.462 (1996) 563. arXiv:arXiv:astro-ph/9508025, doi:10.1086/177173.
* [5] Fermi-LAT collaboration, Constraining Dark Matter Models from a Combined Analysis of Milky Way Satellites with the Fermi Large Area Telescope, Physical Review Letters 107 (24) (2011) 241302. arXiv:1108.3546, doi:10.1103/PhysRevLett.107.241302.
* [6] F. Donato, D. Maurin, P. Brun, T. Delahaye, P. Salati, Constraints on WIMP Dark Matter from the High Energy PAMELA $\overline{p}/p$ Data, Physical Review Letters 102 (7) (2009) 071301. arXiv:0810.5292, doi:10.1103/PhysRevLett.102.071301.
|
arxiv-papers
| 2013-09-18T12:50:26 |
2024-09-04T02:49:51.096290
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Roberto A. Lineros",
"submitter": "Roberto Lineros",
"url": "https://arxiv.org/abs/1309.4630"
}
|
1309.4665
|
# Generalizing the self-healing diffusion Monte Carlo approach to finite
temperature: a path for the optimization of low-energy many-body bases
Fernando A. Reboredo Jeongnim Kim Materials Science and Technology Division,
Oak Ridge National Laboratory, Oak Ridge, TN 37831, USA
###### Abstract
A statistical method is derived for the calculation of thermodynamic
properties of many-body systems at low temperatures. This method is based on
the self-healing diffusion Monte Carlo method for complex functions [F. A.
Reboredo J. Chem. Phys. 136, 204101 (2012)] and some ideas of the correlation
function Monte Carlo approach [D. M. Ceperley and B. Bernu, J. Chem. Phys. 89,
6316 (1988)]. In order to allow the evolution in imaginary time to describe
the density matrix, we remove the fixed-node restriction using complex
antisymmetric guiding wave functions. In the process we obtain a parallel
algorithm that optimizes a small subspace of the many-body Hilbert space to
have maximum overlap with the subspace spanned by the lowest-energy
eigenstates of a many-body Hamiltonian. We show in a model system that the
partition function is progressively maximized within this subspace. We show
that the subspace spanned by the small basis systematically converges towards
the subspace spanned by the lowest energy eigenstates. Possible applications
of this method to calculate the thermodynamic properties of many-body systems
near the ground state are discussed. The resulting basis can be also used to
accelerate the calculation of the ground or excited states with Quantum Monte
Carlo.
###### pacs:
02.70.Ss,02.70.Tt
## I Introduction
There is a significant interest in thermodynamical properties observed as
$T\rightarrow 0$. Many physical phenomena that cover superconductivity,
magnetic and structural transitions, chemical reactions etc. require an
adequate treatment of thermal effects. These effects are crucial in systems
where there is a large number of low-energy excitations within an energy
window $1/\beta=k_{B}T$ above the ground state. Electronic thermal effects are
expected to be larger in metals and magnets than in insulators. QTS In metals
there is a significant number of excitations with vanishing energy. The
magnetic excitations energies frequently go to zero in the long wave limit. A
significant fraction of spectroscopic techniques probe the electronic or
magnetic excitations near the ground state. The development of first-
principles techniques to obtain excitations has historically received a
significant theoretical attention. runge84 ; onida02 ; filippi09 Monte Carlo
methods used to calculate excitation energiesfilippi09 will be accelerated
with basis that retain the physics at the relevant energies.
A first-principles finite-temperature description of many-body systems is also
relevant to describe chemical reactions.mazzola12 Ionic dynamics are usually
calculated within the Born-Oppenheimer approximation. This decouples the wave
function of the “quantum” electrons from the wave function of the ions. Within
this standard approximation, electrons are at zero temperature while the ions
can move with kinetic energies that often exceed the electronic excitations.
Even within the Born-Oppenheimer-ground-state approximation, the standard
approach based on density functional theory (DFT) shows significant
differences with diffusion Monte Carlo (DMC) grossman97 ; saccani13 or
Quantum Chemistry benchmarks. At the transition saddle points, when some
chemical bonds are broken and new ones are formed, the spacing of the
corresponding electronic eigenenergies is minimum, or even zero at the conical
intersections. levine07 Electronic thermal effects are seldom included in
many-body calculations. morales10 ; mazzola12 In order to routinely include
thermal effects, significant improvements in the theory beyond the standard
approach are required.
Most ab-initio calculations in the literature of condensed matter electronic
structure are based in the ground state quantum Monte Carlo calculations of
the homogeneous electron gasceperley80 which made possible the first
approximations of DFT. hohenberg ; kohn DFT has been extended to finite
temperature long ego. mermin65 ; kohn83 ; weinert92 Fermi occupations of
Kohn-Sham eigenstates and the addition of an entropy term have been
shownweinert92 to provide a variational density functional. However, even
nowadays, the zero temperature approximation for the exchange-correlation
potential is widely used. This approach has been long known to severely
underestimate the critical Curie temperatures of magnetic systems. statunton84
; gyorffy85 ; staunton85 Including temperature for magnetic systems is
possible for cases where the magnetic excitations can be treated classically
wang95 ; landau04 and the electrons can be assumed to be in the ground state
for constrained configurations of the spins. eisenbach09 But an adequate
description of the electronic entropy in the subspace that preserves the spin
is still lacking. yin12 Finite temperatures benchmarks of a quality
comparable to Ref. ceperley80, are the key ingredients required to
parametrize a finite temperature density functional. Without a reliable
approximation, most work done under a DFT framework still uses a zero
temperature approximation for the exchange correlation functional.
Accurate many-body calculations at high temperatures can be performed within
the path integral Monte Carlo approach (PIMC). ceperley95 Since the cost of
PIMC diverges as $T\rightarrow 0$, it has been mainly used in the hot and
dense regime, magro96 ; militzer00 ; driver12 with a temperature comparable
to the interaction potential. An alternative approach that could start from
the zero temperature limit would be desirable.
The most accurate techniques to describe a large number of electrons
($N_{e}>1000$) at zero temperature are based in projection approaches.
ceperley80 ; purwanto09 ; booth12 One could potentially extend these methods
to finite temperature, limiting the projector $e^{-\beta\mathcal{\hat{H}}}$ to
finite $\beta$. Thermodynamical averages can be later obtained from
derivatives of the Helmholtz free energy
$F(\beta)=-1/\beta\ln[tr(e^{-\beta\mathcal{\hat{H}}})]$, where
$\mathcal{\hat{H}}$ is the many-body Hamiltonian operator and $tr(X)$ the
trace of $X$ over the complete many-body Hilbert space.
The standard diffusion Monte Carlo Method with importance sampling (DMC)
ceperley80 ; mfoulkesrmp2001 ; austin12 constrains the sign or the phase of
the wave function by imposing the nodes or the phase ortiz93 of a guiding
wave function $\Psi_{T}({\bf R})$, where the many-body coordinate ${\bf
R}=\\{{\bf r}_{1},{\bf r}_{2},\cdots,{\bf r}_{N_{e}}\\}$ is the set of
coordinates of $N_{e}$ electrons. These constraints while enforcing an anti-
symmetric fermionic wave function introduce a variational error. The quality
of the wave function and its nodes can be improved with several methods within
a variational Monte Carlo (VMC) context, mfoulkesrmp2001 ; ortiz95 ; jones97 ;
guclu05 ; umrigar07 ; toulouse08 ; petruzielo12 ; rios06 or at the DMC level.
ceperley88 ; jones97 ; keystone ; rollingstones ; stonehenge In DMC, the
energy of the ground state is exact if the exact nodes or phase are provided.
anderson79 ; ortiz93 Improving the nodes is computationally intensive.
Avoiding this cost is key for finite temperature calculations.
In standard DMC calculations, $\mathcal{\hat{H}}$ is replaced by the fixed-
node Hamiltonian $\mathcal{\hat{H}}_{FN}$ or the fixed-phase Hamiltonian
$\mathcal{\hat{H}}_{FP}$. The use of the fixed-node or fixed phase
approximation can have undesired effects on the calculation of thermal
effects. It has been found that many fermionic systems have a ground state
with two nodal pockets. bressanini12 That is, if the ground state wave
function is real the nodal surface separates the Hilbert space in only two
pockets for positive and negative values respectively. It has been conjectured
ceperley91 that this is a general property of fermionic ground states. In the
fixed-node case, the excitations of $\mathcal{\hat{H}}_{FN}$ are forced to
share the nodes of the ground state. To be orthogonal to the fixed-node ground
state, the fixed-node excited states have to have at least an additional node.
It is known, however, that in many systems there are several fermionic excited
states near the ground state with also two nodal pockets.rasch12 Accordingly,
$tr(e^{-\beta{\mathcal{\hat{H}}}_{FN}})$ does not describe the low temperature
physics. It is easy to see that the same happens in the fixed-phase case.
Therefore, if one wishes to use a DMC-like algorithm to obtain thermodynamical
properties, one must go beyond the usual fixed-node or fixed-phase
approximations. For practical reasons, a parallel approach that can handle a
large number of excitations near the ground state would also be beneficial.
In this paper, we restart the debate on how to calculate low temperature
properties within a many-body ab-initio context taking into account recent
theoretical developments. keystone ; rockandroll ; stonehenge A method is
derived that introduces temperature within an importance sampling procedure
that shares most of the computational tools developed for projection MC
approaches. The errors in the evolution operator $e^{-\beta\mathcal{\hat{H}}}$
resulting from the fixed-node restriction are eliminated by using complex
linear combinations of eigenstates, which do not have nodes except at the
electronic coincidental points (see Fig 1). Instead of optimizing a single
many-body wave-function so it better describes the ground state of the system,
a basis of several wave functions is optimized to maximize the overlap with
the small subspace spanned by the lowest energy eigenstates of the many-body
Hamiltonian. We show in a model system that the overlap of the optimized
subspace with the lowest energy subspace calculated with a configuration
interaction (CI) approach, increases systematically as the iterations increase
and that the partition function is maximized.
The rest of the paper is organized as follows: In Section II we describe the
general formalism; some of the formulae developed in Ref. stonehenge, for
complex wave function is repeated here for completeness. Section III outlines
the basic algorithm. In Section IV we describe the results for a model
calculation; and finally, in Section V we discuss the possible applications
and summarize. This paper also has three appendices: A describes how to go
beyond the locality and local-time approximations; B describes how to take
advantage of the eigenstates when they are complex; Finally, C describes how
to work with eigenstate pairs to minimize the variance of the weights of the
walkers while keeping the wave function complex.
## II A low-energy expansion of the partition function
This section extends the DMC approach ceperley80 for the calculation of the
partition function of a many-body system. We first provide background material
required to understand the rest of the paper. We generalize the upper bound
property of the energy in DMC to an upper bound property of the free energy.
We next give the general outline of our approach and describe how to avoid the
fixed-node approximation in DMC. Finally, we describe the details: basic
formulae and numerical approach.
### II.1 The upper bound property of the truncated Helmholtz free energy
Thermal effects can be obtained by calculating all excitations within a
thermal energy window above the ground state larger than $\Delta E\sim
1/\beta=k_{B}T$ and then evaluating the density matrix reichl as:
$\displaystyle\hat{\rho}(\beta)=$ $\displaystyle e^{-\beta\mathcal{\hat{H}}}$
(1) $\displaystyle=$ $\displaystyle\sum_{n}\left|\Psi_{n}\right\rangle
e^{-\beta E_{n}}\left\langle\Psi_{n}\right|,$
where$E_{n}<E_{0}+\Delta E$ is the eigenvalue with eigenvector
$|\Psi_{n}\rangle$ of $\mathcal{\hat{H}}$. In general $\mathcal{\hat{H}}$ is
given by
$\hat{\mathcal{H}}=\sum_{j}^{N_{e}}\frac{(\nabla_{j}+{\bf
A}_{j})^{2}}{2}+\hat{V}({\bf R})$ (2)
where ${\bf A}_{j}={\bf A}({\bf r}_{j})$ is a vector potential at point ${\bf
r}_{j}$ with magnetic field ${\bf B}({\bf r}_{j})=\nabla_{j}\times{\bf
A}_{j}$, and $\hat{V}({\bf R})$ includes the electron-electron interaction,
the interactions of the electrons spins with the magnetic field and any
external potential, local or non-local.
In a closed system that can exchange energy with a bath or reservoir
(canonical ensemble) all thermodynamical averages can be obtained using the
density matrix. The trace of the density matrix
$Z(\beta)=tr[\hat{\rho(\beta)}]$ is the partition function, whereas
$F(\beta)=-1/\beta\ln[Z(\beta)]$ is the Helmholtz free energy.
In general, $\mathcal{\hat{H}}$ has an infinite number of eigenvectors
$|\Psi_{n}\rangle$ that can be ordered with increasing eigenenergy $E_{n}$. If
$\beta(E_{n}-E_{0})\gg 1$, the contribution to $Z(\beta)$ of the eigenstate
$|\Psi_{n}\rangle$ becomes negligible. Therefore, a usual approximation is to
truncate the trace to a finite matrix with a finite number of eigenstates
$M_{S}$.
In what follows we defined $tr()$ as the trace of a truncated square matrix
with size $M_{S}$. We also relate $Z(\beta)$ and $F(\beta)$ to that truncated
trace. Since $e^{-\beta\mathcal{\hat{H}}}$ is positive definite, for a given
basis, $Z(\beta)$ increases and $F(\beta)$ decreases as $M_{S}$ increases.
The trace of any linear operator is invariant for linear transformations of
the form ${\hat{B}}e^{-\beta\mathcal{\hat{H}}}{\hat{B}}^{-1}$ with
${\hat{B}}{\hat{B}}^{-1}=1$. Thus, in principle, one does not need to obtain
the eigenstates of $e^{-\beta\mathcal{\hat{H}}}$ or equivalently
$\mathcal{\hat{H}}$ to calculate the free energy. Any linearly independent
basis that spans the same subspace can be used to obtain $Z(\beta)$.
Thermodynamical properties only require us to evaluate $Z(\beta)$ in a
linearly independent basis $\\{|\chi^{S}_{m}\rangle\\}$. However, if
statistical methods are used, then each element contributing to the truncated
trace also increases the statistical error bar. Therefore, it is
computationally more efficient to use the most compact basis, with minimum
$M_{S}$, that retains the low-energy properties.
Any eigenstate $|\Psi_{n}\rangle$ can be written in a complete basis
${|\chi^{S}_{n}}\rangle$ as
$\displaystyle|\Psi_{n}\rangle=e^{-\hat{J}}\left[\sum_{m=0}^{M_{S}-1}\lambda_{n}^{m}|\chi^{S}_{m}\rangle+\sum_{m=M_{S}}^{N_{B}\rightarrow\infty}\lambda_{n}^{m}|\chi^{S}_{m}\rangle\right],$
(3)
where $\langle{\bf R}|e^{-\hat{J}}|{\bf R}\rangle=e^{-J({\bf R})}$ is a
Jastrow factor that introduces adequate cusp conditions kato1957 and are
$\langle{\bf R}|\chi^{S}_{m}\rangle=\chi_{m}^{S}({\bf R})$ linear combinations
of an infinite orthogonal set $\\{\langle{\bf R}|n\rangle\\}$ [e.g. Slater
determinants, or Pfaffians bajdich08 , or symmetry constrained functions
(SCF), etc]. In practice we restrict the Hilbert space to a finite number
$N_{B}$. We denote the subspace spanned by $N_{B}\gg M_{S}$ functions
$|m\rangle$ as the big subspace. The big subspace has to be large enough to
describe the low temperature physics of the complete Hilbert space, which is
in general infinite. We define the small subspace as the subspace spanned by
the first $M_{S}$ basis functions $|\chi^{S}_{n}\rangle$.
Within the small basis $\\{|\chi^{S}_{m}\rangle\\}$, the free energy will be
minimum if all $|\Psi_{n<M_{S}}\rangle$ can be spanned in the small basis,
namely, $\lambda_{n\geq M_{S}}^{m}=\lambda_{n}^{m\geq M_{S}}=0$. Errors in the
small basis will result in higher values of the free energy. Thus the free
energy in the truncated basis is an upper bound to the true Helmholtz free
energy. Optimization of the Helmholtz free energy in the small basis is
analogous to the variational principle of the ground state. Likewise, the
partition function in the truncated basis is a lower bound of the exact
partition function. Improved bounds may be obtained with a basis that better
describes the lower energy eigenstates. $M_{S}$ has to be large enough to
include all the relevant physics for a given temperature.
Most of the optimization methods in the QMC literature focus, on optimizing
the eigenstates. Several methods have been proposed to obtain low-energy
excited states within the linear method Monte Carlofilippi09 ; toulouse12
(LMMC) or diffusion Monte Carlo. ceperley88 ; rockandroll However, since the
eigenstates are sometimes difficult calculate, and we only need an average, we
argue that one might save computational time by optimizing the many-body basis
first as in the correlation function diffusion Monte Carlo (CFDMC) ceperley88
method. Optimizing the basis directly could be more practical than obtaining
accurate eigenstates energies, if the number of excitations near the ground
state is large (e. g. typically the case in metallic or magnetic systems).
### II.2 Guiding ideas of the finite temperature SHDMC method and definitions
Instead of performing the usual projection for infinite imaginary time of a
single trial wave-function, we run DMC for multiple guiding wave-functions
(forming a linearly independent basis) for finite imaginary time, which is
equivalent to finite temperature. Instead of using a single real guiding
function with nodes, we use a set of complex antisymmetric guiding functions
without nodes. Therefore, the Hamiltonian $\mathcal{\hat{H}}$ is not altered
at the nodes as in the standard importance sampling DMC approachceperley80
with the fixed-node approximation.anderson79 As explained in the
introduction, extending DMC to finite temperatures requires to go beyond those
fixed schemes. Complex-valued antisymmetric wave functions, that do not have
nodal pockets, can be constructed as a linear combination of two real wave
function with different nodes (see Fig. 1). We go beyond the standard fixed-
phase approximation and the local-time approximation. stonehenge As in the
SHDMC method for complex wave functions, stonehenge the infamous sign problem
is avoided with complex antisymmetric guiding functions. The result is acurate
as long as enough statistical information is collected.
SHDMCkeystone ; rockandroll ; stonehenge systematically improves a trial wave
function by maximizing the overlap with the wave function propagated in
imaginary time in DMC. Here instead of maximizing the overlap of a single wave
function we will maximize the overlap with the basis.
Following the ideas of the SHDMC method, we use a recursive approach. In every
iteration $\ell$, importance sampling DMC ceperley80 is performed and
statistical data of the evolution in $\beta$ of a set of $M_{S}$ guiding wave
functions is projected on the many-body bases
$\\{|\chi^{S,\ell}_{n}\rangle\\}$ and $\\{|n\rangle\\}$ . The statistical data
is used to improve the small basis $\\{|\chi^{S,\ell+1}_{n}\rangle\\}$ and the
guiding functions for the next iteration In what follows, we will omit the
iteration index $\ell$ in the notation for clarity, when the basis is not
changed. The small basis $\\{|\chi^{S}_{n}\rangle\\}$ is orthonormal:
$\langle\chi^{S}_{n}|\chi^{S}_{n}\rangle=\delta_{n,m}$.
For numerical efficiency, depending on the problem, we choose guiding
functions related to the eigenstates of
$\mathcal{\hat{U}}=e^{\hat{J}}e^{\mathcal{\hat{H}}}e^{-\hat{J}}$. (i)
$\\{|\chi^{U}_{n}\rangle\\}$ is formed by the Slater expansions of the
eigenstates in the small basis. (ii) $\\{|\chi^{V}_{n}\rangle\\}$ is formed by
the Slater expansion of linear combinations of eigenstates pairs.
We construct wave functions of the form
$\displaystyle\Psi^{T}_{n}({\bf R})=$ $\displaystyle e^{-J({\bf
R})}\chi^{X}_{n}({\bf R});$ (4)
where the super index $X$ refers to either $S$, $U$, or $V$, depending on the
case. To simplify the notation, we omit $X$ in $\Psi^{T}_{n}({\bf R})$.
We assume the Jastrow factor operator $e^{-\hat{J}}$ to be diagonal in the
many-body configuration space ${\bf R}$, and positive, which implies that it
must have an inverse. fn:back_flow The Jastrow factor is fixed in SHDMC but
can be optimized variationaly so that the free energy of the system is
minimized.
In contrast with the CFDMC ceperley88 and the released phase jones97
methods, we use anti-symmetric guiding functions, which are improved
recursively with a maximum overlap criterion. Since the exponential growth of
the bosonic ground state is prevented by the guiding functions, the free
energy obtained is an upper bound. This approach is different to the
correlated linear method filippi09 because the wave function is optimized at
the DMC level and we use anti-symmetic guiding functions. In variance with the
original SHDMC approach for excited states rockandroll multiple wave
functions are propagated in parallel. A serial orthogonalization step in the
original SHDMC method for excited states rockandroll ; stonehenge is
postponed in this new approach until DMC has been run for the all basis
functions. fn:complications
### II.3 Working with complex guiding wave functions to avoid the fixed-node
approximation
Figure 1: (Color online) A schematic representation of the nodes of the real
and imaginary parts of a complex antisymmetric function. This type of complex
functions must be zero only at the coincidental points. Though real
antisymmetric functions must have nodal pockets a complex linear combination
must not. Note also that two different antisymmetric real wave functions
approaching different eigenstates will have different nodes and produce
different fixed-node Hamiltonians.
While complex guiding wave functions allow us to avoid the fixed-node
approximation, they introduce additional complications stonehenge that are
discussed here. Once these complications are dealt with, the sign problem is
avoided as in the fixed-node. As in standard SHDMC the result is accurate as
long as enough statistics is obtained.
Following Refs. ortiz93, and stonehenge, , $\Psi^{T}_{n}({\bf R})=\langle{\bf
R}|\Psi_{n}^{T}\rangle$ can be written as an explicit product of a complex
phase and an amplitude $\Psi^{T}_{n}({\bf R})=\Phi_{n}^{T}({\bf R})e^{{\bf
i}\phi_{n}({\bf R})}$.
The expressions
$\displaystyle\Phi_{n}^{T}({\bf R})$ $\displaystyle=\sqrt{\Psi_{n}^{T}({\bf
R})\Psi^{T*}_{n}({\bf R})},\text{ and }$ (5) $\displaystyle\phi_{n}({\bf R})$
$\displaystyle=\ln[\Psi_{n}^{T}({\bf R})/\Psi_{n}^{T*}({\bf R})]/(2i)+\pi m$
(6)
allow the computation of all the gradients and Laplacians in terms of those of
an arbitrary complex function $\chi_{n}^{X}({\bf R})$ and $e^{-J({\bf R})}$.
In Eq. (6) $m$ is an arbitrary integer that changes the Riemann branch of the
natural logarithm $\ln$ of a complex number. $m$ only contributes to the
gradient or Laplacian at the Reimann cuts. fn:reimann Since the position of
the Reimann cuts is an arbitrary mathematical convention, their contribution
to gradients and Laplacians is unphysical and ignored.
The dependence in $\beta$ of $e^{-\beta\mathcal{\hat{H}}}\Psi_{n}^{T}({\bf
R})$ is given by
$\displaystyle\Psi_{n}^{T}({\bf R},\beta)=$
$\displaystyle\;e^{-\beta\hat{\mathcal{H}}}\Psi_{n}^{T}({\bf R})$ (7)
$\displaystyle=$
$\displaystyle\;e^{-\beta\hat{\mathcal{H}}}\left[\Phi_{n}^{T}({\bf R})e^{{\bf
i}\phi_{n}({\bf R})}\right]$ (8) $\displaystyle=$
$\displaystyle\;\Phi_{n}^{T}({\bf R},\beta)e^{{\bf i}\phi_{n}({\bf R})}.$ (9)
Equation (9) includes by definition all the temperature dependence in
$\Phi_{n}^{T}({\bf R},\beta)~{}=~{}\Psi_{n}^{T}({\bf R},\beta)/e^{{\bf
i}\phi_{n}({\bf R})}$, denoted as free-amplitude, since it can be complex,
stonehenge with $\Phi_{n}^{T}({\bf R},0)=\Phi_{n}^{T}({\bf R})$. The phase
$\phi_{n}({\bf R})$ remains fixed in the interval $[0,\beta)$ as in Ref.
ortiz93,
Following Ref. stonehenge, we define the quantity
$\displaystyle f_{n}({\bf R},\beta)=\Psi_{n}^{T*}({\bf R})\Psi^{T}_{n}({\bf
R,\beta})e^{\beta E^{X}_{n}},$ (10)
where $E^{X}_{n}$ is a reference energy adjusted numerically to satisfy the
condition
$\langle\chi^{X}_{n}|e^{\hat{J}}|\Psi^{T}_{n}({\beta})\rangle~{}=~{}1$. This
reference energy is different from the one commonly used to obtain the ground
state. In practice, $E^{X}_{n}$ depends on the Slater determinant expansion
$\chi^{X}_{n}({\bf R})$ used to construct the guiding wave function and
contains the relevant information required to calculate thermodynamical
averages.
Using Eqs (2) and (7), one can easily obtain
$\displaystyle\frac{\partial}{\partial\beta}f_{n}({\bf R},\beta)=$
$\displaystyle\Psi_{n}^{T*}({\bf
R})\frac{\partial}{\partial\beta}\left[\Psi^{T}_{n}({\bf R,\beta})e^{\beta
E^{X}_{n}}\right]$ $\displaystyle=$ $\displaystyle-\Psi_{n}^{T*}({\bf
R})\left[(\hat{\mathcal{H}}-E^{X}_{n})\Psi^{T}_{n}({\bf R,\beta})e^{\beta
E^{X}_{n}}\right]$ $\displaystyle=$ $\displaystyle-\left[E_{n}^{T}({\bf
R},\beta)-E^{X}_{n}\right]f_{n}({\bf R},\beta)$ (11)
with
$\displaystyle E_{n}^{T}({\bf R},\beta)$
$\displaystyle=\frac{\mathcal{\hat{H}}\Psi^{T}_{n}({\bf
R},\beta)}{\Psi^{T}_{n}({\bf R},\beta)}$
$\displaystyle=-\frac{1}{2}\\!\\!\sum_{j}^{N_{e}}\frac{\nabla_{j}^{2}\Phi_{n}^{T}({\bf
R},\beta)}{\Phi_{n}^{T}({\bf R},\beta)}$ (12)
$\displaystyle+\frac{1}{2}\sum_{j}^{N_{e}}\left|{\bf
A}_{j}+{\bf\nabla}_{j}\phi_{n}({\bf R})\right|^{2}+V({\bf R},\beta)$
$\displaystyle-{\bf
i}\sum_{j}^{N_{e}}\left\\{\frac{{\bf\nabla}_{j}\Phi_{n}^{T}({\bf
R},\beta)}{\Phi_{n}^{T}({\bf R},\beta)}\cdot\left[{\bf
A}_{j}\\!+\\!{\bf\nabla}_{j}\phi_{n}({\bf R})\right]\right.$
$\displaystyle\;\;\;\;\;\;\;\;\;\;\;\;\;+\left.\frac{{\bf\nabla}_{j}\cdot\left[{\bf
A}_{j}+{\bf\nabla_{j}}\phi_{n}({\bf R})\right]}{2}\right\\}.$
In order to perform an importance sampling using $\Phi_{n}^{T}({\bf R})$ as a
guiding wave function (as in Refs. ceperley80, and ortiz93, ), it is
convenient to express the kinetic energy in terms of $\Phi_{n}^{T}({\bf R})$.
The term including ${\nabla_{j}^{2}\Phi_{n}^{T}({\bf
R},\beta)}/{\Phi_{n}^{T}({\bf R},\beta)}$ can be rewritten using the
identityceperley80 ; HLRbook ; mfoulkesrmp2001
$\displaystyle\frac{\nabla^{2}_{j}\Phi^{T}_{n}({\bf
R},\beta)}{\Phi_{n}^{T}({\bf R},\beta)}=$
$\displaystyle\frac{\nabla^{2}_{j}f_{n}({\bf R,\beta})}{f_{n}({\bf
R,\beta})}+\frac{\nabla^{2}_{j}\Phi^{T}_{n}({\bf R})}{\Phi^{T}_{n}({\bf R})}$
$\displaystyle-\frac{{\bf\nabla}_{j}\cdot\left[f_{n}({\bf R},\beta){\bf
F}_{Q}^{j}\right]}{f_{n}({\bf R,\beta})}\;,$ (13)
where
$\displaystyle{\bf F}_{Q}^{j}={\bf\nabla}_{j}ln\left|\Phi_{n}^{T}({\bf
R})\right|^{2}\;.$ (14)
Note that Eqs. (II.3) and (14) are valid as long as $\Phi^{T}_{n}({\bf R})\neq
0$ and $f_{n}({\bf R},\beta)\neq 0$. In practice, any divergence of ${\bf
F}_{Q}^{j}$ at the nodes enforces $f_{n}({{\bf R},\beta})$ to be zero. Figure
1 shows that complex antisymmetric wave functions can be constructed so that
they have nodes only at points ${\bf R}$ with ${\bf r}_{i}={\bf r}_{j}$. In
this case the nodal error is avoided but errors in the phase introduce a phase
shiftortiz93 ; stonehenge and a complex contribution to $E^{L}_{n}({\bf
R},\beta)$. However, for complex wave functions without zeros, Eq. (II.3) is
always valid, except at the coincidental points (if cusp conditions are not
satisfied). To satisfy Eq. (II.3) at the coincidental points, a Jastrow factor
is introduced in Eq. (4). While using complex wave functions involves some
complications, the advantage is that the evolution in imaginary time $\beta$
describes the thermodynamical properties with $\beta^{-1}=k_{B}T$. However,
going beyond the fixed-phase approximation ortiz93 ; meyer13 is required to
obtain the thermodynamics. In this work the phase is not “released” in the
same sense of Ref. jones97 , it is only free within the small subspace.
Replacing Eq. (II.3) into Eq (II.3) and then into Eq. (II.3) one obtains
$\displaystyle\frac{\partial f({\bf R},\beta)}{\partial\beta}=$
$\displaystyle\frac{1}{2}\sum_{j}^{N_{e}}\left\\{\nabla_{j}^{2}f_{n}({\bf
R},\beta)-{\bf\nabla}_{j}\cdot\left[f_{n}({\bf R},\beta){\bf
F}_{Q}^{j}\right]\right\\}$ $\displaystyle-\left[E_{n}^{L}({\bf
R},\beta)-E^{X}_{n}\right]f_{n}({\bf R},\beta),$ (15)
with $E_{n}^{L}({\bf R},\beta)$ where
$\displaystyle E_{n}^{L}({\bf R},\beta)=$
$\displaystyle-\frac{1}{2}\\!\\!\sum_{j}^{N_{e}}\frac{\nabla_{j}^{2}\Phi_{n}^{T}({\bf
R})}{\Phi_{n}^{T}({\bf R})}$ (16)
$\displaystyle+\frac{1}{2}\sum_{j}^{N_{e}}\left|{\bf
A}_{j}+{\bf\nabla}_{j}\phi_{n}({\bf R})\right|^{2}+V({\bf R},\beta)$
$\displaystyle-{\bf
i}\sum_{j}^{N_{e}}\left\\{\frac{{\bf\nabla}_{j}\Phi_{n}^{T}({\bf
R},\beta)}{\Phi_{n}^{T}({\bf R},\beta)}\cdot\left[{\bf
A}_{j}\\!+\\!{\bf\nabla}_{j}\phi_{n}({\bf R})\right]\right.$
$\displaystyle\;\;\;\;\;\;\;\;\;\;\;\;\;+\left.\frac{{\bf\nabla}_{j}\cdot\left[{\bf
A}_{j}+{\bf\nabla_{j}}\phi_{n}({\bf R})\right]}{2}\right\\}.$
That is, the local energy now depends on two guiding functions (i)
$\Phi_{n}^{T}({\bf R})$, and (ii) $\Phi_{n}^{T}({\bf R},\beta)$ which is an
approximation that must be obtained and improved for $\beta\neq 0$.
The use of complex valued guiding functions originates the gradients
${\bf\nabla}_{j}\phi_{n}({\bf R})$ that appear in the local energy in Eq.
(16). Their contribution prevents the result from reaching the bosonic
solution and enforces an upper bound on the fermionic ground state. ortiz93
In addition, the contribution of ${\bf\nabla}_{j}\Phi_{n}^{T}({\bf R},\beta)$
must be taken into account in the presence of a magnetic field (see term
between the $\\{\\}$, when ${\bf A}_{j}\neq 0$), even when using a real-
bosonic guiding wave function with ${\bf\nabla}_{j}\phi_{n}({\bf R})=0$.
A locality approximation locality has been used in the past when a non-local
pseudo potential is included in $\mathcal{\hat{H}}$ in the potential term
$\hat{V}({\bf R})$. It consists in replacing $V({\bf
R},\beta)=\frac{\hat{V}({\bf R})\Psi^{T}_{n}({\bf R},\beta)}{\Psi^{T}_{n}({\bf
R},\beta)}$ by $V({\bf R},0)$ in Eq.(II.3). Since $V({\bf R},\beta)\rightarrow
V({\bf R},0)$ for $\beta\rightarrow 0$, we will use the locality approximation
in the first iteration. However, we will improve it in subsequent iterations
(see Appendix A).
A local-time approximation, analogous to the locality approximation, locality
, was introduced in Ref. stonehenge, to estimate the ratio
$\displaystyle\frac{{\bf\nabla}_{j}\Phi_{n}^{T}({\bf
R},\beta)}{\Phi_{n}^{T}({\bf R},\beta)}=$
$\displaystyle\frac{{\bf\nabla}_{j}\Phi_{n}^{T}({\bf R})}{\Phi_{n}^{T}({\bf
R})}+{\bf\nabla}_{j}\ln\left[\frac{\Phi_{n}^{T}({\bf
R},\beta)}{\Phi_{n}^{T}({\bf R})}\right].$ (17)
One can neglect the last term in Eq. (17) for $\beta\rightarrow 0$ where
${\Phi_{n}^{T}({\bf R},\beta)}/{\Phi_{n}^{T}({\bf R})}\rightarrow 1$. In the
present work we will use the local-time approximation only in the first
iteration. For subsequent iterations we improve the evaluation of Eq. (17) by
using the sampling of the dependence in $\beta$ of $\Phi_{n}^{T}({\bf
R},\beta)$ obtained in the previous iteration.
The locality and local-time approximations have little impact in optimization
methods that focus on eigenstates because the dependence on $\beta$ of
$\Psi^{T}_{n}({\bf R,\beta})$ is minimized when the optimization progresses as
$\Psi^{T}_{n}({\bf R})\rightarrow\Psi_{n}({\bf R})$. Going beyond these
approximations is required, however, to circumvent the nodes of the
eigenstates with complex wave functions. Fortunately, optimization in the
small subspace allows an easy sampling of the $\beta$ dependence. The approach
is exact if the big basis is large enough and if enough statistical data is
collected fn:fixedphase as $\ell\rightarrow\infty$.
Circumventing the nodes with complex wave functions is necessary in this case
because, in standard DMC calculations using real-valued wave functions with
nodes, if any walker crosses a node, it is either killedceperley80 or the
move is rejected. umrigar93 This introduces an artificial divergent potential
at the nodal surface, which adds a kink at the node (a step for the rejection
case). Since there is a one-to-one correspondence between energy of one
eigenstate and its nodes, rosetta eigenstates with different energies must
have different nodes. As a consequence, two real wave wave functions that
approach different eigenstates introduce different nodal potentials. Since the
fixed-node Hamiltonian is different for different eigenstates, and affect the
dynamics at the node in the evolution in imaginary time, the $\beta$
dependence obtained using the fixed-node approximation will not describe the
thermodynamics even if the exact nodes of the ground state are provided.
### II.4 Differences with other DMC-like projection methods
The implementation of this method follows essentially the same approach
developed for DMC or SHDMC, with some key numerical changes.
Equation (15) is identical to Eq. (1) in Ref. ceperley80, except for the
local energy, which now has an explicit dependence in $\beta$. Unlike Eq. (13)
in Ref. stonehenge, , Eq. (15) is now valid for $\beta\geq 0$. As in Ref
ceperley80, , Eq (15) describes the evolution of an ensemble. Each member of
the ensemble of configurations ${\bf R}_{i}$ (walker) undergoes (i) a random
diffusion and (ii) drifting by the quantum force $\sum_{j}{\bf F}_{Q}^{j}$
(which depends only on $\Phi_{T}({\bf R})$ and not on the phase). Following
Ref. stonehenge, , (iii) each walker carries a complex phase factor. In a
nonbranching algorithm, the complex weight of the walkers is multiplied by
$\exp\\{-\left[E^{L}_{n}({\bf R},\beta)-E^{X}_{n}\right]\delta\beta\\}$ at
every diffusion step. In contrast with Ref. jones97, , the phase factor of the
walkers starts in 1 and evolves towards the difference between the guiding
phase and the phase of $\Psi({\bf R},\beta)$, while in the release phase
approach the initial phase of each walker depends on the initial positions of
the walkers ${\bf R}_{i}$ but remains constant in $\beta$.
If the $\chi_{n}^{X}({\bf R})$ are linear combinations of antisymmetric
functions, with arbitrary complex coefficients, the $\Phi_{n}^{T}({\bf R})$ in
Eq. (16) do not have nodal surfaces (see Fig. 1). Therefore, ${\bf F}_{q}$ is
not divergent, but at the coincidental points.
All walkers must add the same inverse temperature $\beta$ after $k$ steps. The
standard time-step correction to minimize time step errors at the nodes [Eq.
(33) in Ref. umrigar93, ] is not strictly necessary since the divergences in
${\bf F}_{q}$ are removed. If one introduces it, one must readjust
$\delta\beta$ during the time evolution so that all the walkers add up to the
same $\beta$. For the same reason, the standard accept/reject scheme that
enforces detailed balanceHLRbook is modified: other moves are retried after
rejection until a move is accepted.
### II.5 Calculation of the partition function
This section shows that the partition function can be obtained as
$Z(\beta)\simeq\sum_{n=0}^{M_{S}-1}e^{-\beta E^{S}_{n}}$ where the $E^{S}_{n}$
are reference energies instead of eigen energies.
Numerically, it is convenient to start the calculation with a distribution of
walkers proportional to $f_{n}({\bf R},0)=|\Phi^{T}_{n}({\bf R})|^{2}$. In the
importance sampling approach ceperley80 setting the second line of Eq. (15)
equal to zero, provides an equilibrium distribution proportional to
$|\Phi^{T}_{n}({\bf R})|^{2}$.
As in the DMC and SHDMC methods, the evolution in inverse temperature $\beta$
is discretized into $k$ finite steps $\delta\beta=\beta/k$. Following the
SHDMC approachkeystone ; rockandroll ; rollingstones ; stonehenge the
weighted distribution of the walkers can be written as
$\displaystyle f_{n}({\bf R},k\delta\beta)$
$\displaystyle=\lim_{N_{c}\rightarrow\infty}\frac{1}{N_{c}}\sum_{i=1}^{N_{c}}W_{i}\delta\left({\bf
R-R}_{i}\right).$ (18)
In Eq. (18), ${\bf R}_{i}$ corresponds to the position of the walker $i$ , and
$N_{c}$ is the number of equilibrated configurations. The complex weights
$W_{i}$ are given by
$\displaystyle W_{i}=e^{-\left[E_{i}^{k}-E^{X}_{n}\right]k\delta\beta}$ (19)
with
$\displaystyle E_{i}^{k}=\frac{1}{k}\\!\sum_{j=0}^{k-1}E^{L}_{n}({\bf
R}_{i}^{-j}),$ (20)
Where $k$ is a number of steps and $E_{L}({\bf R}_{i}^{-j})$ is the previous
value of the local energy obtained $j$ steps $\delta\beta$ earlier for the
walker $i$.
The evolution in inverse temperature $\beta$ of the guiding wave function
$\Psi_{n}^{T}({\bf R},\beta)$ can be written, without loss of generality, as
$\displaystyle\Psi_{n}^{T}({\bf R},\beta)=e^{-\beta E^{X}_{n}}e^{-J({\bf
R})}\left[\chi^{X}_{n}({\bf R})+\delta\chi^{X}_{n}({\bf R},\beta)\right].$
(21)
That is, the product of an average decay factor $e^{-\beta E^{X}_{n}}$ times
the Slater determinant part. The Slater part is given by the the one at
$\beta=0$ plus an orthogonal displacement $\delta\chi^{X}_{n}({\bf R},\beta)$.
The $X$ in term $\delta\chi^{X}_{n}({\bf R},\beta)$ denotes the explicit
dependence on $\chi^{X}_{n}({\bf R})$. At least one overlap
$\langle\chi^{X}_{m}|\delta\chi^{X}_{n},\beta\rangle$ must be different from
zero for $n\neq m$ , if $\Psi^{T}_{n}({\bf R})$ is not an eigenstate.
Using Eq. (21) we can correct equation (16) beyond the locality and local-time
approximations. The displacement $\delta\chi^{X}_{n}({\bf R},\beta)$ can be
sampled from the DMC run as follows: From Eqs. (7), (10) and (18), one can
obtain
$\displaystyle e^{-J({\bf R})}\delta\chi^{X}_{n}({\bf R},\beta)=$
$\displaystyle\langle{\bf R}|\left[e^{\beta
E^{X}_{n}}|\Psi^{T}_{n}(\beta)\rangle-|\Psi^{T}_{n}\rangle\right]$
$\displaystyle=$ $\displaystyle\;\langle{\bf
R}|\left[e^{-\beta(\hat{\mathcal{H}}-E^{X}_{n})}-1\right]e^{-\hat{J}}|\chi^{X}_{n}\rangle\;$
(22) $\displaystyle=$ $\displaystyle\;\frac{\left[f_{n}({\bf
R},\beta)-f_{n}({\bf R},0)\right]}{\Psi^{T*}_{n}({\bf R})}$ (23)
$\displaystyle\simeq$ $\displaystyle\frac{1}{N_{c}}\sum_{i=1}^{N_{c}}e^{J({\bf
R})}\frac{[W_{i}-1]}{\chi^{X*}_{n}({\bf R})}\delta({\bf R}-{\bf R}_{i}).$ (24)
Within the subspace spanned by the basis $\\{\chi^{S}_{n}({\bf R})\\}$, the
identity operator $\hat{E}$ is given by
$\displaystyle\langle{\bf R}^{\prime}|\hat{E}|{\bf R}\rangle$
$\displaystyle=\sum_{m=0}^{M_{S}-1}e^{-J({\bf R^{\prime}})}\chi^{S}_{m}({\bf
R^{\prime}})\chi^{S*}_{m}({\bf R})e^{J({\bf R})}.$ (25)
Applying Eq. (25) to both sides of Eq. (24), and integrating over ${\bf R}$,
one can easily obtain an expression of the diffusion displacement within the
basis $\\{|\chi^{S}_{n}\rangle\\}$ as
$\displaystyle\delta\chi^{X}_{n}({\bf R},\beta)$
$\displaystyle=\sum_{m=0}^{M_{S}-1}\lambda_{n}^{m}(\beta)\chi^{S}_{m}({\bf
R})$ (26)
with
$\displaystyle\lambda_{n}^{m}(\beta)=\frac{1}{N_{c}}\sum_{i=1}^{N_{c}}e^{2J({\bf
R}_{i})}\frac{\chi^{S*}_{m}({\bf R}_{i})}{\chi^{X*}_{n}({\bf
R}_{i})}[W_{i}-1]$
where $N_{c}=\sum_{i=1}^{N_{c}}W_{i}$.
Using $X=S$ in Eqs. (18) -(21) one can prove that
$\displaystyle\langle\chi^{S}_{n}|\mathcal{\hat{U}}|\chi^{S}_{m}\rangle=$
$\displaystyle\int{\bf dR}e^{J({\bf R})}\chi^{S*}_{n}({\bf
R})\Psi_{m}^{T}({\bf R},\beta)$ $\displaystyle=$ $\displaystyle e^{-\beta
E^{S}_{n}}\left[\delta_{n,m}+\int{\bf dR}\chi^{S*}_{n}({\bf
R})\delta\chi^{S}_{m}({\bf R},\beta)\right]$ $\displaystyle=$ $\displaystyle
e^{-\beta E^{S}_{n}}(\delta_{n,m}+\lambda_{n}^{m})$ (28)
with $\mathcal{\hat{U}}$ having the structure of the transcorrelated method
ten-no02
$\displaystyle\mathcal{\hat{U}}=$ $\displaystyle
e^{\hat{J}}e^{-\beta\mathcal{\hat{H}}}e^{-\hat{J}}.$ (29)
We use condition $\lambda_{n}^{n}=0$ [See Eq (LABEL:eq:deltalambda)] to
determine the value of $E^{S}_{n}$. In practice, we adjust the reference
energy of the guiding functions every iteration as $e^{-\beta
E^{X,\ell+1}_{n}}=e^{-\beta E^{X,\ell}_{n}}(\delta_{n,m}+\lambda_{n}^{m})$.
Since $tr(e^{-\beta\mathcal{\hat{H}}})=tr(\mathcal{\hat{U}})$, the
contribution to the Helmholtz free energy of the small subspace is given by
$\displaystyle
F(\beta)=-\frac{1}{\beta}\ln{\left[\sum_{n=0}^{M_{S}-1}e^{-\beta
E^{S}_{n}}\right]},$ (30)
where the expression inside the brackets is the partition function $Z(\beta)$.
In general for an arbitrary guiding function $\Psi_{m}^{T}({\bf R})$, the
variance will grow with $M_{S}$. An energy span larger than $\beta^{-1}$ must
be retained in the basis to calculate thermodynamical properties. Arbitrary
trial wave functions spanned by this space might have significant variance in
the walkers weights. To reduce the variance we use guiding functions that are
approximately linear combinations of a pair of neighboring eigenstates.
When using guiding functions that are different from the small basis
functions, the contributions to the trace of the density matrix in the small
basis can be obtained with
$\displaystyle e^{-\beta
E^{S}_{n}}=\sum_{i}\left[|\langle\chi^{X}_{i}|\chi^{S}_{n}\rangle|^{2}+\langle\chi^{X}_{i}|\chi^{S}_{n}\rangle\langle\chi^{S}_{n}|\delta\chi^{X}_{i}\rangle\right]e^{-\beta
E^{X}_{i}}\;.$ (31)
The details of the derivation are in Appendix B.
In a recent paper, Mazzola, Zen and Sorellamazzola12 proved that
$\displaystyle\langle\Psi^{T}_{n}|e^{-\beta\hat{\mathcal{H}}}|\Psi^{T}_{n}\rangle\geq
e^{-\beta\langle\Psi^{T}_{n}|\hat{\mathcal{H}}|\Psi^{T}_{n},\rangle}.$ (32)
Ref. mazzola12, used the righthand side of Eq. (32) to approximate the free
energy obtaining a lower bound for $F(\beta)$. Reference mazzola12, can be
considered a VMC approach to the evaluation of the free energy. That
approximation becomes exact if all the $|\Psi^{T}_{n}\rangle$ are eigenstates
of $\hat{\mathcal{H}}$. However, that method is very poor for an arbitrary
random guiding function. In the present approach, we go beyond Ref. mazzola12,
by evaluating the lefthand side of Eq. (32) directly using DMC.
In many situations, the excitations of a mean field method based on
approximations DFT might be good enough to obtain the low energy
thermodynamical properties using Eq. (30). If that were the case, at least two
DMC runs for each function of the basis are required. One to obtain the
$\beta$ dependence and a second to evaluate the reference energies beyond the
local-time approximation. However, in the so-called highly correlated
materials, usual approximations of DFT fail to describe the low energy
physics. In those cases a method that could optimize the basis is more
important. That method is described in the following subsections.
### II.6 The first iteration: Construction of the first small basis
$\\{|\chi^{S,1}_{n}\rangle\\}$
While the present approach will optimize the basis from any starting basis
set, the calculation will be more efficient starting from a good basis. A
procedure to generate a good starting set is described here.
The only restriction for the small basis $\\{|\chi^{S}_{n}\rangle\\}$ is to
avoid the nodes associated with real wave functions. In this work we choose
the initial basis with a Lanczos-like procedure combined with the SHDMC
approach.
The big subspace basis set ${|m\rangle}$ is constructed by symmetry
constrained functions (linear combinations of Slater determinants with the
same symmetry of the ground state) ordered with increasing mean field energy.
We choose the first basis function of the small subspace to be
$\displaystyle|\Psi_{0}^{S}\rangle=\frac{1}{\sqrt{2}}(|0\rangle+{\bf
i}|1\rangle)$ (33)
being $|0\rangle$ and $|1\rangle$the ground and first excited states of a non-
interacting solution of the system.
Using Eq. (33) as guiding function in Eqs. (II.5)-(24) and replacing
$\chi^{S}_{m}({\bf R})$ by $\langle{\bf R}|m\rangle$ in Eq. (25), but not on
$\Psi^{T}_{n}({\bf R})$, we obtain an expression similar to Eq. (26)
$\displaystyle\delta\tilde{\chi}_{n}^{S}({\bf R})=\langle{\bf
R}|\delta\tilde{\chi}^{S}_{n}\rangle=\tilde{\sum}_{m=0}c_{n}^{m}\langle{\bf
R}|m\rangle.$ (34)
The tilde in $\delta\tilde{\chi}_{n}^{S}({\bf R})$ means that the expansion is
in the big basis $\\{|m\rangle\\}$ with
$\displaystyle c_{n}^{m}(\beta)=\frac{1}{N_{c}}\sum_{i=1}^{N_{c}}e^{2J({\bf
R}_{i})}\frac{\langle m|{\bf R}_{i}\rangle}{\chi^{S*}_{n}({\bf
R}_{i})}[W_{i}-1].$
The symbol $\tilde{\sum}$ in Eq. (34) means that the sum is restricted to the
coefficients $c_{n}^{m}$ with an error bar smaller than $25\%$ of the absolute
value (this is the standard recipe of the SHDMC algorithm keystone ).
We define the next basis function $|\chi^{S}_{n+1}\rangle$ recursively as
$\displaystyle|\chi^{S}_{n+1}\rangle=$
$\displaystyle\frac{1}{\mathcal{N}_{n+1}}\hat{P}_{n+1}|\delta\tilde{\chi}^{S}_{n}\rangle\text{
with }$ (36) $\displaystyle\hat{P}_{n}=$ $\displaystyle
1-\sum_{m=0}^{n-1}|\chi^{S}_{m}\rangle\langle\chi^{S}_{m}|$ (37)
where $\mathcal{N}_{n+1}$ is a normalization constant. Equations (37) and (36)
mean that $|\chi^{S}_{n}\rangle$ is the projection of the displacement
$|\delta\tilde{\chi}^{S}_{n+1}\rangle$ orthogonal to the subspace spanned by
the $n$ basis functions found previously.
One repeats this procedure until a basis of $M_{S}$ functions
$\\{|\chi^{S}_{n}\rangle\\}$ is constructed. This Lanczos-like procedure
grants that the initial small subspace basis times the Jastrow factor has a
large projection onto the lowest-energy eigenstates of $\mathcal{\hat{H}}$ or
the largest eigenstates of $e^{-\beta\mathcal{\hat{H}}}$.
Since the evolution in inverse temperature $\beta$ is not known during the
initialization step, we use the local-time approximation discussed in the
previous section. However, once a basis is generated, we can go beyond the
local-time approximation in successive iterations. Note that by construction
any $\chi^{S}_{n}({\bf R},\beta)$ can be approximated as a linear combination
of the basis function $\chi^{S}_{m}({\bf R})$ with $m<n+1$, since the finite
temperature projection of one wave function of the basis into the other was
used to construct the small basis. The details on how to approximate the
evolution in $\beta$ are in Appendix A.
### II.7 Systematic improvement of the small basis
$\\{|\chi^{S}_{n}\rangle\\}$
One of the goals of this work is to obtain a much smaller basis of $M_{S}$
functions $\\{|\chi^{S}_{n}\rangle\\}$ than $\\{|m\rangle\\}$, the big set of
$N_{B}$ basis functions. The small basis $\\{|\chi^{S}_{n}\rangle\\}$ should
retain the lowest energy physics of $\mathcal{\hat{H}}$ and
$e^{-\beta\mathcal{\hat{H}}}$. While for some purposes (e.g. the ground state
calculation), the initial basis set described in the previous section might be
enough, in this subsection we describe how to further optimize the small basis
so it better describes the low-energy excitations of $\mathcal{\hat{H}}$ and
thermodynamical properties.
Note that $\chi^{S}_{n}({\bf
R},\beta)=e^{\hat{J}}e^{-\beta\mathcal{\hat{H}}}e^{-\hat{J}}\chi^{S}_{n}({\bf
R})$ will converge to the antisymmetric part of the ground state wave-function
as $\beta\rightarrow\infty$. In order to avoid every state in the basis
collapsing to the same function we (i) remove the projection into the other
states of the basis, (ii) add the diffusion displacement orthogonal to the
small subspace, and (iii) perform a Gram Schmidt orthogonalization as follows:
$\displaystyle|\chi^{S,\ell+1}_{0}\rangle=$
$\displaystyle|\chi^{S,\ell}_{0}\rangle+|\delta\tilde{\chi}_{0}^{S}\rangle-|\delta\chi^{S}_{0}(\beta)\rangle$
(38) $\displaystyle|\chi^{S,\ell+1}_{n}\rangle=$
$\displaystyle\hat{P}^{\ell+1}_{n}\left[|\chi^{S,\ell}_{n}\rangle+|\delta\tilde{\chi}_{n}^{S}\rangle-|\delta\chi^{S}_{1}(\beta)\rangle\right]\;\;n>0,$
with $\hat{P}^{\ell+1}_{n}$ given by Eq. (37) replacing $\chi^{S}_{m}$ by
$\chi^{S,\ell+1}_{m}$. Note in Eq. (38) that
$|\delta\tilde{\chi}_{n}^{S}\rangle$ given by Eq. (34) is a direct projection
of the diffusion displacement into the big basis of SCFs $\\{|m\rangle\\}$,
whereas $|\delta\chi^{S}_{n}(\beta)\rangle$ given by Eq. (26) is an indirect
projection (since the $|\delta\chi^{S}_{n}(\beta)\rangle$ are projected into
the small basis $\\{|\chi^{S}_{n}\rangle\\}$ which in turn are linear
expansions of functions that belong to $\\{|m\rangle\\}$). The difference
$|\Delta\tilde{\chi}^{S}_{n}\rangle=|\delta\tilde{\chi}^{S}_{n}\rangle-|\delta\chi^{S}_{n}(\beta)\rangle$
is by construction orthogonal to the small subspace. Accordingly, it describes
the decay of the small subspace basis into the eigenstates with the lowests
energies.
## III Algorithm
The goal of this algorithm is to optimize a minimal basis
$\\{|\chi^{S}_{n}\rangle\\}$ to span the lowest energy excitations of
$\mathcal{\hat{H}}$ (equivalently the eigenstates of
$e^{-\beta\mathcal{\hat{H}}}$ with the largest eigenvalues). That basis can be
used to calculate finite temperature expectation values of thermodynamical
properties and accelerate the calculation of the ground and lower excited
states. In this section we summarize how the theory described in detail
earlier can be implemented.
Initialization: the small basis $\\{|\chi^{S}_{n}\rangle\\}$, a set of
orthogonal linear combinations of many-body functions $|m\rangle$ is
constructed using the procedure described in II.6. Once this procedure is
concluded, a set of linearly independent guiding functions
$e^{-\hat{J}}|\chi^{V}_{n}\rangle$ is constructed as linear combinations of
pairs of approximated eigenstates of $\mathcal{\hat{U}}$.
Subsequently, we can use for the evaluation of the local energy
$E^{L}_{n}({\bf R},\beta)$ an approximate dependence of the guiding functions
in $\beta$ given by Eq. (47).
Basis update iteration: each iteration $\ell$ is composed by (1) a parallel
diffusion of intermediate functions $|\chi^{V}_{n}\rangle$. (2) A linear
transformation to obtain $|\delta\tilde{\chi}^{S}_{n}\rangle$ and
$|\delta\chi^{S}_{n}\rangle$. (3) Recalculation of $\mathcal{\hat{U}}$ in the
small basis. (4) Update of the small basis for the next iteration
$|\tilde{\chi}^{S,\ell+1}_{n}\rangle$. (5) Update of the intermediate
functions $|\chi^{V,\ell+1}_{n}\rangle$. (6) Finally, the algorithm decides to
increase the number of samples in the next iteration or not. These six steps
are repeated recursively.
The individual steps of the iteration are described below in more detail:
(1) Parallel diffusion: Each displacement
$\displaystyle|\delta\chi^{V}_{n}\rangle=e^{\hat{J}}\left[e^{-\beta\mathcal{\hat{H}}}-1\right]e^{-\hat{J}}|\chi^{V}_{n}\rangle$
(39)
is projected into the small basis $\\{|\chi^{S}_{n}\rangle\\}$ and the big
basis $\\{|m\rangle\\}$ using Eqs. (LABEL:eq:deltalambda) and Eq.
(LABEL:eq:cnm) respectively replacing $\chi^{X}_{n}({\bf R})$ by
$\chi^{V}_{n}({\bf R})$.
Each diffusion contains $M_{b}$ sampling subblocks. For each sampling
subblock, uncorrelated walker positions are generated from the previous one
with a VMC algorithm. Next, DMC is run for $k$ steps with a shorter time step
$\delta\beta$. The coefficients of $|\delta\tilde{\chi}^{V}_{n}\rangle$ and
$|\delta\chi^{V}_{n}\rangle$ are sampled at the end of each subblock using
Eqs. (LABEL:eq:deltalambda) and (LABEL:eq:cnm). Statistical data is collected
for $M_{b}$ subblocks for each parallel diffusion before an update of the
small basis.
(2) Linear transformation: The $|\delta\tilde{\chi}^{S}_{n}\rangle$ and
$|\delta\chi^{S}_{n}\rangle$ can be constructed in terms of
$|\delta\tilde{\chi}^{V}_{n}\rangle$ and $|\delta\chi^{V}_{n}\rangle$ using
Eq. (B) replacing the super index $X$ by $V$.
(3) Calculation of $\mathcal{\hat{U}}$, $\\{|\Psi^{U,\ell+1}_{n}\rangle\\}$,
and $\\{|\Psi^{V,\ell+1}_{n}\rangle\\}$ : A matrix representation of
$\mathcal{\hat{U}}$ is obtained in the small subspace
$\\{|\Psi^{U,\ell}_{n}\rangle\\}$, using Eqs. (LABEL:eq:deltalambda), (II.5)
and (B). The left and right eigenvectors are obtained by diagonalizing
$\mathcal{\hat{U}}$.
(4) Update of the small Basis: Equation (38) is used to improve the small
basis.
(5) We perform the correspondence
$\chi^{S,\ell}_{n}\rightarrow\chi^{S,\ell+1}_{n}$ (see Appendix A ) and
construct the basis $\\{|\chi^{U,\ell+1}_{n}\rangle\\}$, and
$\\{|\chi^{V,\ell+1}_{n}\rangle\\}$ with the coefficients of the eigenvectors
of $\mathcal{\hat{U}}$ in the iteration $\ell$.
(6) Updating $M_{b}$: At first, the number of sampling subblocks $M_{b}$ is
set to a small number (e.g. $M_{b}=3$). When the noise is dominant
$\sum_{n}\langle\Delta\tilde{\chi}^{S,\ell+1}_{n}|\Delta\tilde{\chi}^{S,\ell}_{n}\rangle\leq
0$, we increase $M_{b}$ by a factor larger than 1. As a result, the total
number of configurations $N_{c}$ sampled increases as the iteration $\ell$
increases and the statistical error is reduced. Hence, as the statistics are
improved the number of basis functions retained in the expansion Eq. (25)
increases over time.
(7) We use Eq. (30) to calculate thermodynamical properties.
## IV Results in model calculations
This section describes the results obtained for a model system with an applied
magnetic field. The results are compared with configuration interaction (CI)
calculations in the same model used in Refs. rosetta, ; keystone, ;
rockandroll, and stonehenge, . The lowest-energy eigenstates were found for
two polarized electrons ($J=1$) moving in a two-dimensional square with a side
length $1$ and a repulsive interaction potential of the form $\hat{V}({\bf
r},{\bf
r^{\prime}})=8\pi^{2}\gamma\cos{[\alpha\pi(x-x^{\prime})]}\cos{[\alpha\pi(y-y^{\prime})]}$
with $\alpha=1/\pi$ and $\gamma=4$.
The main advantage of the model is that fully converged CI calculations can be
performed which are nearly analytical. In order to perform the CI calculations
the many-body wave function of the small basis $\\{\chi^{S}_{n}\rangle\\}$ are
spanned in a basis of functions $\\{|m\rangle\\}$ that are eigenstates of the
noninteracting system. They are linear combinations of functions of the form
$\prod_{\nu}\sin(m_{\nu}\pi x_{\nu})$ with $m_{\nu}\leq 7$. Converged CI
calculations were performed to obtain a nearly exact expression of the lowest
energy states of the system $|\Psi_{n}\rangle=\sum_{m}a_{m}^{n}|m\rangle$. The
matrix elements involving the magnetic vector potential ${\bf A}$ (in the
symmetric gauge) were calculated analytically. The result of the CI
calculations were used to evaluate the partition functions and to quantify the
convergence of the basis.
The same basis used to construct the CI Hamiltonian is used as the big basis
to test our finite temperature version of SHDMC. All the calculations reported
are with $J({\bf R})=0$, which increases the statistically noise, makes the
test more difficult and facilitates the comparison with the CI results. The
results presented here are a proof of principle on the validity of the
algorithm, which is necessary before requesting and using the massive amount
of computing time required for realistic finite temperature calculations in
solids. While clearly a demonstration in a realistic system is required in the
future, a comparison with an exact model is the first essential step to
validate the scheme. This includes not only the value obtained for the
partition function but also a detailed analysis of the convergence of the
basis.
In the absence of magnetic fields there are two degenerate solutions: one that
transforms line $x$, and the other that transform like $y$. This degeneracy is
broken with a magnetic field. The eigenstates transform like $x+{\bf i}y$ and
$x-{\bf i}y$. Figure 2 shows the evolution of
$Z_{+}(\beta)=tr_{+}(\mathcal{\hat{U}})$ of the model system with $M_{S}=20$.
The subindex “$+$” in $Z_{+}$ and $tr_{+}$ means that the results of Fig. (2)
were obtained considering only the subspace of the Hamiltonian that transforms
like $x+{\bf i}y$. The calculation of thermodynamical properties requires,
however, the inclusion of all possible symmetries of the wave function, which
implies that $Z_{-}$, the trace in a small basis that transforms as $x-{\bf
i}y$, should also be added. In order to calculate $Z_{+}(\beta)$ we have
defined the zero of energy to be the ground state of the CI. The calculations
were run using $\delta\beta=0.00002$ and $\beta=0.004$. We have used a
magnetic field of $B=0.6283$.
Figure 2 shows the value of $Z_{+}(\beta)$, calculated with different methods,
relative to the full CI value obtained with the $M_{S}=20$ lowest eigenvalues.
The blue cycles were obtained with SHDMC using Eq. (30). The red rhombi were
obtained by evaluating
$Z_{+}(\beta)=\sum_{n,m}|\langle\Psi^{T}_{m}|\Psi_{n}\rangle|^{2}e^{-\beta
E_{n}}$, being the $|\Psi_{n}\rangle$ and $E_{n}$ the eigenvectors and
eigenvalues of the full CI. The empty squares were obtained as
$Z_{+}(\beta)\simeq\sum_{n}e^{-\beta\langle|\Psi^{T}_{n}|\hat{\mathcal{H}}|\Psi^{T}_{n}\rangle}$
being the $|\Psi^{T}_{n}\rangle$ linear combinations of pairs of approximated
eigenstates obtained with SHDMC. Therefore, the squares correspond to the
result that one would had obtained using the approximation of Ref. mazzola12,
in Eq. (32) for a very good set of functions. Finally, the up triangles mark
the result obtained with
$Z_{+}(\beta)\simeq\sum_{n}e^{-\beta\langle|\Psi^{T}_{n}|\hat{R^{-1}}\hat{\mathcal{H}}\hat{R}|\Psi^{T}_{n}\rangle}$,
being $\hat{R}$ a random rotation defined in the small subspace.
Figure 2 shows that as the iteration $\ell$ increases, the $Z_{+}(\beta)$
obtained with all methods increases. Note, that the scale of the $y$ axis does
not start from zero. The initialization scheme already produces a basis that
retains 90 % of the exact value of the truncated partition function. Similar
to previous SHDMC methods, the present generalization optimizes the small
basis overlap, not their average energy. The trace of $\mathcal{\hat{U}}$
increases indirectly as the small basis approaches to the subspace of the
eigenvectors with lowest energy. Note that the SHDMC results are within $\sim
1.5$% below the values obtained analytically by projection into the CI data
(red rhombi). This difference is due to the remaining errors in the evolution
phase which neglects the projection orthogonal to the small subspace. The
method used in Ref. mazzola12, applied to approximated pairs of eigenstates
(empty squares) gives results only slightly below the SHDMC values, because
the energy difference between eigenstates is much smaller that $\beta^{-1}$.
However, a random rotation of the basis that spans the same subspace (up
triangles) would had produced a significantly worse result. This shows that
the approach described in Ref. 12 is significantly worse if each element of
the basis does not have a large projection on each eigenstate.
Note in the inset of Fig. 2 that the number of sampling blocks remains very
low ($M_{b}=3$) for the first 25 iterations and starts increasing when the
noise becomes dominant around $\ell=32$. If one considers only the first 32
iterations a significant improvement of the partition function is achieved
with little computational cost as compared with the calculation of individual
eigenstates.
Figure 2: (Color online) Evolution of the trace of
$\mathcal{\hat{U}}=e^{\hat{J}}e^{-\beta\mathcal{\hat{H}}}e^{-\hat{J}}$ in the
subspace defined by the small basis $\\{|\chi^{S,\ell}_{n}\rangle\\}$ as a
function of the iteration index $\ell$ for $M_{S}=20$ relative to the
truncated trace obtained with full CI eigenvalues. Note that the scale starts
from 86% of the full CI value. The blue circles are results obtained with
SHDMC using Eq. (30). The red rhombi denote the exact evaluation of
$\mathcal{\hat{U}}=e^{\hat{J}}e^{-\beta\mathcal{\hat{H}}}e^{-\hat{J}}$ in the
small basis, using full CI data. The empty squares (up triangles) were
obtained by evaluating
$\langle\Psi^{T}_{n}|\hat{\mathcal{H}}|\Psi^{T}_{n},\rangle$ for linear
combination of eigenstate pairs (random rotation within the small subspace) in
the full CI basis while $Z(\beta)$ was obtained using an lower bound
approximation. mazzola12 The inset shows the computational cost, the number
of sampling subblocks $M_{b}$, as a function of the iteration.
Figure 3 quantifies the convergence of the small basis
$\\{|\chi^{S,\ell}_{n}\rangle\\}$ towards the basis defined by the eigenstates
of the CI $\\{|\Psi_{n}\rangle\\}$. For that purpose we define the logarithm
of the residual subspace projection as
$\displaystyle
LR_{sp}=\ln\left\\{1-|Det\left[\langle\Psi_{n}|\chi^{S,\ell}_{m}\rangle\right]|^{1/M_{S}}\right\\}.$
(40)
In Eq. (40) $Det\left[\langle\Psi_{n}|\chi^{S,\ell}_{m}\rangle\right]$ is the
Determinant of a square matrix of size $M_{S}$ formed by the overlap
$\langle\Psi_{n}|\chi^{S,\ell}_{m}\rangle$. The determinant of the matrix is a
complex number of modulus $1$ in the limit when any eigenstate
$|\Psi_{n<M_{S}}\rangle$ can be written as a linear combination of
$|\chi^{S,\ell}_{m}\rangle$. Any error in any member of the small basis
$\\{|\chi^{S,\ell}_{m}\rangle\\}$ reduces the modulus of the determinant by a
factor. The exponent $1/M_{S}$ in Eq. (40) is a standard geometric average. A
large negative value in Eq. (40) indicates a very good small basis with a
determinant that is approaching $1$.
Figure 3 shows the evolution of $LR_{sp}$ given by Eq. (40) as a function of
the iteration index $\ell$ for the same system described in Fig. 2. Note that
$LR_{sp}$ becomes increasingly negative as a function of $\ell$, which implies
a global improvement of the basis approaching to the one described by the
eigenstates of the full CI.
Figure 3: Logarithm of the residual subspace projection $LR_{sp}$ [see Eq.
40] of the small basis $\\{|\chi^{S,\ell}_{n}\rangle\\}$ into the basis
$\\{|\Psi_{n}\rangle\\}$ defined by the eigenstates of the full CI as a
function of the iteration index $\ell$ for the same system described in Fig 2.
A larger negative number correspond to in improved small basis.
We next need to characterize how well an individual eigenstate can be
described by the small basis. To measure this we define the logarithm of the
residual projection $LR_{n}$ as
$\displaystyle
LR_{n}=\ln\left[1-\sum_{m=0}^{M_{S}-1}|\langle\Psi_{n}|\chi^{S,\ell}_{m}\rangle|^{2}\right].$
(41)
Note in Eq. (41) that, if the normalized eigenstate can be written as a linear
combination of the small basis $\\{|\chi^{S,\ell}_{m}\rangle\\}$, the
expression in the brackets should be zero. A large negative number in $LR_{n}$
implies that the eigenstate $|\Psi_{n}\rangle$ is very well described in the
small basis.
Figure 4 describes the evolution of $LR_{n}$ for different eigenstates of the
CI as a function of the iteration index $\ell$.
The blue (red) contribution to the color decreases (increases) as the index
$n$ increases. The continuous line follows the ground state. One can clearly
observe that the ground state of the CI is already very well described at the
initialization stage within the Lanczos-like setup. As the iteration $\ell$
increases, the small basis describes the lowest-energy excitations better
while higher excitations require more iterations. Note that for 25 iterations
15 eigenvectors are very well described within a basis of 20 states. The total
cost at this point is 300 000 individual DMC steps. The calculations of 15
eigenstates with the original SHDMC algorithm for excited statesrockandroll
would had cost at least twice as much. The current algorithm becomes
competitive, in addition, if one considers that it is parallel, which allows
to distribute this cost in multiple tasks ($M_{S}$) reducing the time to
solution to 2% as compared with the original SHDMC algorithm for excited
states.
Figure 4: (Color online) Evolution of the logarithm of residual projection
$LR_{n}$ as a function of the iteration index $\ell$. Blue (red) color denotes
lower (higher) values of $n$. The continuous line follows the ground state
eigenstate.
Finally, for infinite statistics one could in principle obtain the
eigenenergies of $\mathcal{\hat{H}}$ from the eigenvalues of $u_{n}$ as
$E_{n}\approx-\ln(u_{n})/\beta$. This procedure is known to be inefficient to
obtain the eigenenergies which are better described by sampling
$\mathcal{\hat{H}}$ as in the CFDMC approach. The off diagonal noise in the
matrix elements of $\mathcal{\hat{U}}$ has a perverse effect on the magnitude
of the eigenvalues. Therefore, while this method is an efficient one to
optimize the basis, it should be combined with other methods to obtain the
eigenvalue spectra.
## V Summary and discussions
In this paper we have presented a general framework aimed to calculate
thermodynamical properties of many-body system in an importance sampling DMC
context. ceperley80
We showed that a many-body basis can be optimized to describe a small subspace
maximizing the overlap with the subspace described by the lowest eigenstates
of the Hamiltonian. The Helmholtz free energy obtained within this truncated
basis is an upper bound of the exact free energy of a system. The
corresponding partition function is a lower bound of the exact partition
function.
This generalization of the SHDMC method for finite temperature takes advantage
of complex wave functions that do not have nodal pockets. Accordingly, we
avoid the appearance of the artificial potentials in the standard fixed-node
approximation when the amplitude of the importance sampling guiding function
is zero. The antisymmetric properties of the wave-function are enforced by a
complex phase. This introduces a complex contribution to the local energy. The
complex local energy is handled in the complex weight of the walkers. Going
beyond the local-time approximation used in Ref. stonehenge, , the evolution
of the complex phase factor for $\beta>0$ is now approximately taken into
account in the evaluation local energy. The evolution in $\beta$ becomes exact
as the statistical error is reduced as the sampling increases.
In variance CFDMC, ceperley88 the solution remains an upper bound of the
fermionic free energy. While the CFDMC approach uses a bosonic-trial wave
function without nodes, SHDMC uses a complex linear combination of anti-
symmetric functions without nodal pockets. The walker distribution is
prevented to fall into the bosonic ground state solution by the phase factor
of the guiding function, which remains antisymmetric, and introduces an
effective potential in the local energy. ortiz93 ; stonehenge
The present approach shares many aspects of the SHDMC method for complex wave
functions, but it also incorporates a key advantage of the CFDMC ceperley88
approach: several wave functions are optimized at the same time. In systems
where many excitations can be approximated by a single amplitude $\Phi({\bf
R})$ and a different phase factor $e^{{\bf i}\phi_{n}({\bf R})}$, a correlated
sampling approach that reweighs the walkers in Eq. (18) as $W_{i}\rightarrow
W_{i}\Phi^{T}_{n}({\bf R})/\Phi({\bf R})$ and changes the phase contribution
to the local energy in Eq. (16), would save significant time. If that
approximation were used, this generalization of the SHDMC method would look
very similar in spirit with CFDMC, the main difference being the use of a
complex guiding functions that prevents the exponential growth of the bosonic
ground state.
The displacement of each wave function in the small basis during the DMC
process is decomposed into a displacement within the subspace already
described by the other elements of the small basis plus a contribution
orthogonal to the small subspace. The displacement included within the basis
is used to improve upon the locality locality and local-timestonehenge
approximations. The displacement orthogonal to the small subspace is used to
correct the small basis used in the next iteration.
The serial orthogonalization step required in the original SHDMC algorithm for
excited states rockandroll ; stonehenge is avoided with a method that allows
the calculation of multiple wave functions in parallel. In addition, the
complications of inequivalent nodal pockets of excited states rockandroll is
avoided using complex trial wave functions without nodes.
The scaling of the cost of an individual iteration of this method is
proportional to $M_{S}\times N_{c}$; $M_{S}$ the size of the small basis and
$N_{c}$ the number of configurations of the DMC run. The cost of and
individual DMC step is dependent of the number of basis functions $N_{b}$ and
electrons.
It is well known that as the size of the system increases, the number of basis
functions $N_{b}$ required to maintain a fixed error bar for a given
eigenstate must increase factorially. But in practice, the error required for
evaluation of thermodynamical properties is determined by $\beta^{-1}$: errors
must be much smaller than the temperature of the system. Therefore, as the
temperature increases and averages of multiple eigenstates are obtained, the
detail required by calculations of the ground state energies with chemical
accuracy is no longer necessary. The present approach can take advantage of
the acceleration of the algorithms used to evaluate large numbers of
determinants. nukala09 ; clark11 For a very large $N_{b}$, the cost of these
algorithms scales as $N_{e}\times N_{b}$.
The total cost is dependent on the physical system and the goal of the
calculation. If the goal is to converge the entire basis or to optimize the
Free energy, the bottleneck for convergence is the energy gap
$E_{M_{S}+1}-E_{M_{S}}$ which determines the convergence of the basis towards
the highest eigenstates considered. Accordingly, in this case, the ideal
situation for this method would be a system with ($M_{S}$) nearly degenerate
eigenstates well-separated from the rest of the spectra within an energy scale
of $k_{b}T$. If the goal, instead, is to converge the small basis so as only
the lowest $M_{L}$ eigenstates are well described, the convergence of the
algorithm is much faster and it is limited by the number of statistical
samples and the exponential decay $e^{-\beta(E_{M_{S}+1}-E_{M_{L}})}$. The
cost is reduced as compared with the calculation of eigenstates if one accepts
an error in the higher excitations. If one wishes to retain the physics of
higher eigenstates in the basis, it is computationally more efficient to
increase $M_{S}$ (which increases the cost linearly), instead of improving the
basis for the higher excitations which increases the cost exponentially.
Comparisons of the method with full CI calculations show that SHDMC can be
used to optimize many-body basis sets to maximize the overlap with the lowest
energy excitations of the Hamiltonian. Each eigenenergy obtained with this
method has lower quality than those obtained with alternative approaches such
as LMMC or the standard SHDMC for excited states. However, this method could
be a useful tool to optimize the basis, minimizing the size of the matrices
used in LMMC and thus reducing the effects of numerical noise in LMMC.
###### Acknowledgements.
The authors would like to thank J. Krogel and P. R. C. Kent for a critical
reading of the manuscript and discussions. This work has been supported by the
grant ERKCS92 Materials Science and Engineering division of Basic Energy
Sciences, Department of Energy.
## Appendix A Going beyond the locality and local-time approximations
Approximate coefficients for the Slater determinant expansion of the
eigenstates of $\mathcal{\hat{H}}$ and $e^{-\beta\mathcal{\hat{H}}}$ can be
obtained from the eigenstates of
$\mathcal{\hat{U}}=e^{\hat{J}}e^{-\beta\mathcal{\hat{H}}}e^{-\hat{J}}$ in the
$\\{\chi^{S}_{n}\\}$ basis [see Eq. (II.5)]. But, $\mathcal{\hat{U}}$ is not
hermitian, since
$\mathcal{\hat{U}}^{\dagger}=e^{2\hat{J}}\mathcal{\hat{U}}e^{-2\hat{J}}$.
Nevertheless, as long as the Jastrow factor operator $e^{-\hat{J}}$ has an
inverse, $\mathcal{\hat{U}}$ has a set of right eigenvectors
$|\chi^{U}_{i}\rangle=e^{\hat{J}}|\Psi_{i}\rangle$ and left eigenvectors
$\langle\chi^{U}_{i}|=\langle\Psi_{i}|e^{-\hat{J}}$. Since $\mathcal{\hat{H}}$
is Hermitian, its eigenstates $|\Psi_{i}\rangle$ are orthogonal, which implies
that $\langle\chi^{U}_{i}|\chi^{U}_{j}\rangle=\delta_{i,j}$. Within
statistical error bars, in the small subspace, the matrix elements of
$\mathcal{\hat{U}}$ obtained with the basis $|\chi^{S,\ell}_{n}\rangle$ of the
iteration $\ell$ are given by Eq. (II.5). In the first iteration the matrix
elements of $\mathcal{\hat{U}}$ can be obtained directly from the Lanczos-like
procedure.
Within the small subspace, $\mathcal{\hat{U}}$ can be written as
$\displaystyle\mathcal{\hat{U}}=\sum_{i}u_{i}|\chi^{U}_{i}\rangle\langle\chi^{U}_{i}|.$
(42)
Since the $u_{i}$ are also the eigenvalues of $e^{-\beta\mathcal{\hat{H}}}$
their dependence with $\beta$ is exponential. Thus for an arbitrary
$\beta^{\prime}$ the eigenvalue will be $u_{i}^{\beta^{\prime}/\beta}$.
Provided that the difference
$|\Delta\tilde{\chi}^{S,\ell+1}_{n}\rangle=|\chi^{S,\ell+1}_{n}\rangle-|\chi^{S,\ell}_{n}\rangle$
is small [which is always valid for $\beta\rightarrow 0$ see Eq. (38)], the
dependence in $\beta$ of $|\chi^{S,\ell+1}_{n}(\beta)\rangle$ can be
approximated as follows: Let’s first define the operator
$\hat{R}=\sum_{n=0}^{M_{S}-1}|\chi^{S,\ell}_{n}\rangle\langle\chi^{S,\ell+1}_{n}|$
and its inverse within the small subspace
$\hat{R}^{\dagger}=\sum_{n=0}^{M_{S}-1}||\chi^{S,\ell+1}_{n}\rangle\langle\chi^{S,\ell}_{n}|$.
Accordingly, the dependence in $\beta^{\prime}$ of the new basis can be
approximated as
$\displaystyle\left|\chi^{S,\ell+1}_{n}\left(\beta^{\prime}\right)\right\rangle=$
$\displaystyle\hat{R}^{\dagger}\mathcal{\hat{U}}^{{\beta^{\prime}}/{\beta}}\hat{R}\left|\chi^{S,\ell+1}_{n}\right\rangle,$
(43)
with $\mathcal{\hat{U}}^{{\beta^{\prime}}/{\beta}}$ given by Eq. (42)
replacing $u_{i}$ by $u_{i}^{\beta^{\prime}/\beta}$.
## Appendix B Working with eigenstates of $\mathcal{\hat{U}}$
While in some systems eigenstates of $\mathcal{\hat{H}}$ are always real (e.g.
confined systems with time reversal symmetry), in many cases the wave function
of the eigenstates is known to be complex, without nodal pockets. In those
cases Eq. (II.3) is valid and no walker needs to be killed or rejected because
the eigenstate wave function does not have a nodal surface. stonehenge In
these cases it might be advantageous to propagate single eigenstates of
$\hat{R}^{\dagger}\mathcal{\hat{U}}^{\frac{j}{k}}\hat{R}$, since the variance
of the weights is minimized and lower (larger) $T$ ($\beta$) can be reached
with less statistical data. An additional advantage of working with functions
that are closer to the eigenstates is that the locality and local-time
approximations can be used.
Since any
$|\chi^{S}_{n}\rangle=\sum_{i=0}^{M_{S}-1}\langle\chi^{X}_{i}|\chi^{S}_{n}\rangle|\chi^{X}_{i}\rangle$
can be written as a linear combination of $|\chi^{X}_{i}\rangle$ and vise
versa, when eigenstates are complex one can use as trial wave function
$e^{\hat{J}}|\chi^{U}_{i}\rangle$ in the SHDMC propagation and sampling. Then
the propagation of $|\chi^{S}_{n}\rangle$ can be written as a linear
combination of the propagation of the eigenstates of $\mathcal{\hat{U}}$. as:
$\displaystyle|\delta\chi^{S}_{n}\rangle$
$\displaystyle=\sum_{i=0}^{M_{S}}\left\\{\langle\chi^{X}_{i}|\chi^{S}_{n}\rangle
e^{-\beta(E^{X}_{i}-E^{S}_{n})}\left[|\chi^{X}_{i}\rangle+|\delta\chi^{X}_{i}\rangle\right]\right\\}-|\chi^{S}_{n}\rangle\text{
and }$ $\displaystyle|\delta\tilde{\chi}^{S}_{n}\rangle$
$\displaystyle=\sum_{i=0}^{M_{S}}\left\\{\langle\chi^{X}_{i}|\chi^{S}_{n}\rangle
e^{-\beta(E^{X}_{i}-E^{S}_{n})}\left[|\chi^{X}_{i}\rangle+|\delta\tilde{\chi}^{X}_{i}\rangle\right]\right\\}-|\chi^{S}_{n}\rangle$
(44)
$E^{S}_{n}$ results from the condition
$\langle\chi^{S}_{n}|\delta\chi^{S}_{n}\rangle=0$ and it is given by
$\displaystyle e^{-\beta
E^{S}_{n}}=\sum_{i}\left[|\langle\chi^{X}_{i}|\chi^{S}_{n}\rangle|^{2}+\langle\chi^{X}_{i}|\chi^{S}_{n}\rangle\langle\chi^{S}_{n}|\delta\chi^{X}_{i}\rangle\right]e^{-\beta
E^{X}_{i}}\;,$ (45)
replacing $X$ by $U$ in Eqs. (B) and (45). The brackets
$\langle\chi^{U}_{i}|\chi^{S}_{n}\rangle$ are obtained by diagonalizing
$\mathcal{\hat{U}}$: the coefficient $n$ of the left eigenvectors $i$ in the
small basis $\\{\langle\chi^{S}_{n}|\\}$.
The disadvantages is that complex eigenstates appear only in certain
Hamiltonians or boundary conditions. Albeit without nodes, they might have
large variation in probability density, in particular for small magnetic
fields or twist boundary conditions close to high symmetry points. Large
variations in the probability density hinder correlating sampling.
## Appendix C Working with eigenstates pairs
It is well known that in many physical systems the energy spacing between
eigenstates decreases as the size of the system increases. It is also known
that the variance of the local energy, which is related to the statistical
error in the energy, increases as the size of the system increases. nemec10
Accordingly as the size of the system increases, it becomes more difficult to
obtain eigenstates. As the size of the system increases the error in the
variance introduced by a linear combination of eigenstates in the Slater part
of the wave function becomes smaller than the variance introduced by short
range correlations. These short range correlations cannot be accounted by the
Slater part, even with a very large basis $\\{|n\rangle\\}$, or with simple
Jastrow factors. Accordingly, in large systems little is lost by using a
linear combination of eigenstates, since their contribution to the variance is
proportional to the energy separation that decreases as the system become
larger. In contrast, much is gained avoiding the divergences at the nodes by
using complex linear combinations of eigenstates, in particular, to obtain
average of thermodynamical properties. However, to propagate for larger
$\beta$ with a branching algorithm, it will be necessary to minimize the
variance of the local energy.
If the eigenstates are real, one must use a linear combination of eigenstates.
The minimum variance will be reached by constructing the Slater part of the
guiding wave function with linear combinations of eigenstates with consecutive
eigenvalues of $u_{i}$. In this work we use
$\displaystyle|\chi^{V}_{2j-1}\rangle$
$\displaystyle=\frac{1}{\sqrt{2}}\left(|\chi^{U}_{2j-1}\rangle+e^{{\bf
i}\phi}|\chi^{U}_{2j}\rangle\right)$ (46) $\displaystyle|\chi^{V}_{2j}\rangle$
$\displaystyle=\frac{1}{\sqrt{2}}\left(|\chi^{U}_{2j-1}\rangle-e^{{\bf
i}\phi}|\chi^{U}_{2j}\rangle\right),$
where $e^{{\bf i}\phi}$ is a complex phase which can be adjusted to minimize
the variance of the amplitude of the complex wave function $\Phi^{T}_{n}({\bf
R})$. The conjugate vectors $\langle\chi^{V}_{2j-1}|$ are constructed using
complex conjugate coefficients and the left eigenvectors of
$\mathcal{\hat{U}}$ in the small basis $\\{\langle\chi^{S}_{n}|\\}$.
Their evolution in imaginary time is given by
$\displaystyle|\chi^{V}_{2j-1}(\beta^{\prime})\rangle$
$\displaystyle=\frac{1}{\sqrt{2}}\left[(u_{2j-1})^{\frac{\beta^{\prime}}{\beta}}|\chi^{U}_{2j-1}\rangle+e^{{\bf
i}\phi}(u_{2j})^{\frac{\beta^{\prime}}{\beta}}|\chi^{U}_{2j}\rangle\right]$
(47) $\displaystyle|\chi^{V}_{2j}(\beta^{\prime})\rangle$
$\displaystyle=\frac{1}{\sqrt{2}}\left[(u_{2j-1})^{\frac{\beta^{\prime}}{\beta}}|\chi^{U}_{2j-1}\rangle-e^{{\bf
i}\phi}(u_{2j})^{\frac{\beta^{\prime}}{\beta}}|\chi^{U}_{2j}\rangle\right]$
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|
arxiv-papers
| 2013-09-18T14:48:38 |
2024-09-04T02:49:51.113414
|
{
"license": "Public Domain",
"authors": "Fernando A. Reboredo and Jeongnim Kim",
"submitter": "Fernando Reboredo",
"url": "https://arxiv.org/abs/1309.4665"
}
|
1309.4741
|
# Overlap Cycles for Permutations:
Necessary and Sufficient Conditions
Victoria Horan
Air Force Research Laboratory
Rome, NY 13441 USA [email protected]
###### Abstract
Universal cycles are generalizations of de Bruijn cycles and Gray codes that
were introduced originally by Chung, Diaconis, and Graham in 1992. They have
been developed by many authors since, for various combinatorial objects such
as strings, subsets, permutations, partitions, vector spaces, and designs. One
generalization of universal cycles, which require almost complete overlap of
consecutive words, is $s$-overlap cycles, which relax such a constraint. In
this paper we study permutations and some closely related class of strings,
namely juggling sequences and functions. We prove the existence of $s$-overlap
cycles for these objects, as they do not always lend themselves to the
universal cycle structure.
## 1 Introduction
††Approved for public release; distribution unlimited: 88ABW-2013-4003, 9 Sep
13
Listing structures for combinatorial objects are quickly becoming useful in
more and more interesting applications. Gray codes, first defined in 1947 by
Frank Gray [6], are used in many different surprising places from position
encoders [11] to genetic algorithms [4]. More recently universal cycles are
being used in areas such as rank modulation for multilevel flash memories
[12], however overlap cycles are still being explored and have the potential
to be useful in many applications.
An $s$-overlap cycle, or $s$-ocycle, is an ordering of a set of objects
$\mathcal{C}$, each represented as a string of length $n$. The ordering
requires that object $b=b_{0}b_{1}\ldots b_{n-1}$ follow object
$a=a_{0}a_{1}\ldots a_{n-1}$ only if $a_{n-s}a_{n-s+1}\ldots
a_{n-1}=b_{0}b_{1}\ldots b_{s-1}$. Ocycles were introduced by Godbole,
Knisley, and Norwood in 2010 [5]. Universal cycles, or ucycles are
$(n-1)$-ocycles and were originally introduced in 1992 by Chung, Diaconis, and
Graham [2].
To find $s$-ocycles and ucycles on a set of strings, most proofs employ the
same method. For a given string $X=x_{1}x_{2}\ldots x_{n}$, let
$X^{s-}=x_{1}x_{2}\ldots x_{s}$ denote the $s$-prefix of $X$ and
$X^{s+}=x_{n-s}x_{n-s+1}\ldots x_{n}$ denote the $s$-suffix of $X$. The first
step in the proof is to construct the transition digraph for the set of
strings as follows. Vertices represent $s$-prefixes and $s$-suffixes of
strings (the overlaps), while each edge represents a string, traveling from
its $s$-prefix to its $s$-suffix. Note that the transition digraph is a
directed multigraph in which an Euler tour (a closed walk in which every edge
is traversed exactly once) corresponds bijectively to an $s$-ocycle. To prove
the existence of an Euler tour, we use the following well-known theorem.
###### Theorem 1.1.
_([13], p. 60)_ A directed graph $G$ is eulerian if and only if it is both
balanced and weakly connected.
In this paper, we consider using the ocycle listing structure for
permutations, as well as functions and juggling sequences. We represent
permutations of an $n$-set $\\{0,1,\ldots,n-1\\}$ as a string
$\Pi=\pi_{0}\pi_{1}\ldots\pi_{n-1}$, where the functional representation is
used, i.e. $\pi(0)=\pi_{0}$. Closely related to permutations, juggling
sequences have been an active research area since the 1980’s [1]. These
sequences are used to determine what patterns a fixed set of balls can be
juggled in, where only one ball may be thrown at a time.
###### Definition 1.2.
_[3]_ A juggling sequence is a string $T=t_{0}t_{1}\ldots t_{n-1}$ where each
$t_{i}\geq 0$ such that
$\left|\\{i+t_{i}\pmod{n}\mid 0\leq i\leq n-1\\}\right|=n.$
This sequence illustrates that at time $i$, we should throw a ball high enough
that it is in the air for $t_{i}$ beats, or to height $t_{i}$. The number of
balls used in a given juggling sequence $T=t_{0}t_{1}\ldots t_{n-1}$ is given
by
$b=\frac{1}{n}\sum_{i=0}^{n-1}t_{i}.$
The period of a juggling sequence $t_{0}t_{1}\ldots t_{n-1}$ is $n$, the
length of the string.
An alternative definition considers the corresponding permutation sequence for
a juggling sequence. Given a string $T=t_{0}t_{1}\ldots t_{n-1}$, the
permutation sequence is the string $\Pi_{T}=\pi_{0}\pi_{1}\ldots\pi_{n-1}$
where $\pi_{i}=t_{i}+i\pmod{n}$. Then $\Pi_{T}$ is a permutation of the
$n$-set $\\{0,1,\ldots,n-1\\}$ if and only if $T$ is a valid juggling
sequence. From this definition it is clear that there is a very close
relationship between juggling sequences and permutations, which is further
illustrated by their similar results explored in this paper. The permutation
sequence clarifies the point that a juggler cannot catch or throw two objects
simultaneously.
For example, when $n=3$ and $b=2$ the pattern $015$ is valid. We can see this
by drawing the juggling diagram, shown in Figure 1. However, the sequence
$105$ is not valid, as shown in Figure 2. In this second example, note that on
the second beat we are required two catch two balls simultaneously - an
operation that is not allowed.
015015015 Figure 1: Juggling Diagram for Sequence $015$ 105105105 Figure 2:
Juggling Diagram for Sequence $105$
In [3], Chung and Graham show that it is not always possible to find a single
universal cycle that contains all juggling sequences of period $n$ and at most
$b$ balls. However they do prove that we can use several disjoint universal
cycles to cover each sequence exactly once. This result is closely related to
the long studied problem of universal cycles for permutations. Since ucycles
for permutations are not possible using the standard representation, a
modification of the representation can often provide an effective solution [9,
10].
An alternative that has only recently been explored is to modify the listing
structure rather than the object representation. By utilizing $s$-ocycles
rather than ucycles, partial results for permutations have been obtained in
[7], however we obtain a more complete solution in this paper. In Section 2 we
establish a complete result on the existence of $s$-ocycles for permutations,
and in Section 3 we explore similar results for functions and juggling
sequences.
## 2 Permutations
We begin with a general lemma that will be used for both permutations and
juggling sequences. For a string $X=x_{0}x_{1}\ldots x_{n-1}$, define the
rotation function as follows. A rotation by $s$ is given by:
$\rho^{s}(X)=x_{s}x_{s+1}\ldots x_{n-1}x_{0}x_{1}\ldots x_{s-1}.$
The following lemma shows that for a string of length $n$, rotations by $s$
partition the string into blocks of length gcd$(n,s)$, and we can always
perform repeated $s$-rotations to start $X$ with any block desired.
###### Lemma 2.1.
Let $n,s\in\mathbb{Z}^{+}$ with $1\leq s\leq n-1$ and gcd$(n,s)=d$. Consider a
string $X=x_{0}x_{1}\ldots x_{n-1}$, also written as $X=Y_{0}Y_{1}\ldots
Y_{m-1}$ where $n=md$ and $Y_{i}=x_{id}x_{id+1}\ldots x_{id+d-1}$. Then for
any $i\in\\{0,1,\ldots,m-1\\}$ there is some $j\in\\{0,1,\ldots,m-1\\}$ such
that:
$Y_{i}Y_{i+1}\ldots Y_{m-1}Y_{0}Y_{1}\ldots Y_{i-1}=\rho^{js}(X).$
###### Proof.
Suppose that $s=kd$. Then each rotation $\rho^{s}(X)$ advances through $k$
blocks of $X$. Thus if we can show that gcd$(m,k)=1$, we are done. However we
note that gcd$(n,s)=d$ implies that there are integers $p,q$ such that
$pn+qs=d$. This implies the following.
$\displaystyle pn+qs$ $\displaystyle=$ $\displaystyle d$ $\displaystyle
pmd+qkd$ $\displaystyle=$ $\displaystyle d$ $\displaystyle pm+qk$
$\displaystyle=$ $\displaystyle 1$
Thus the same integers $p,q$ provide confirmation that gcd$(m,k)=1$.
∎
###### Theorem 2.2.
Let $n,s\in\mathbb{Z}^{+}$ with $1\leq s\leq n-1$ and let $M$ be a multiset of
size $n$. There exists an $s$-ocycle on permutations of $M$ if and only if
$n-s>\hbox{gcd}(n,s)$.
###### Proof.
Define gcd$(n,s)=d$, and suppose that $n-s=kd>d$. Construct the transition
digraph $D$ with vertices representing $s$-permutations of $M$ and edges
representing permutations of $M$. We will show that this graph is balanced and
connected, and hence eulerian. Recall that an Euler tour in this graph
corresponds to an $s$-ocycle for permutations of $M$.
Balanced:
Let $X^{s-}=x_{0}x_{1}\ldots x_{s-1}$ be an arbitrary vertex in the transition
graph. For each possible out-edge corresponding to an $n$-permutation
$X=x_{0}x_{1}\ldots x_{n-1}$, we also have the in-edge
$X^{\prime}=x_{s}x_{s+1}\ldots x_{n-1}x_{0}x_{1}\ldots x_{s-1}$. Thus there is
a one-to-one correspondence between in- and out-edges at each vertex, and so D
is balanced.
Connected:
Consider an arbitrary vertex $X^{s-}=x_{0}x_{1}\ldots x_{s-1}$, which is an
$s$-prefix of some permutation $X=x_{0}x_{1}\ldots x_{n-1}$. We will consider
a partition of $X$ into blocks of size $d$, i.e. $X=X_{0}X_{1}\ldots X_{m-1}$
where $X_{i}=x_{id}x_{id+1}\ldots x_{(i+1)d-1}$. Our goal is to show that we
can permute elements in any $k$ consecutive blocks $X_{i}X_{i+1}\ldots
X_{i+k-1}$, which illustrates that any adjacent transpositions are possible,
and hence we can reach all permutations.
In D, rotations of $X$ correspond to the following cycle, where subscripts are
computed modulo $n$.
$\displaystyle x_{0}x_{1}\ldots x_{s-1}$ $\displaystyle\rightarrow$
$\displaystyle x_{n-s}x_{n-s+1}\ldots x_{n-1}$ $\displaystyle\rightarrow$
$\displaystyle x_{n-2s}x_{n-2s+1}\ldots x_{n-s-1}$ $\displaystyle\rightarrow$
$\displaystyle x_{n-3s}x_{n-3s+1}\ldots x_{n-2s-1}$ $\displaystyle\vdots$
$\displaystyle\rightarrow$ $\displaystyle x_{n-is}x_{n-is+1}\ldots
x_{n-(i-1)s-1}$ $\displaystyle\vdots$ $\displaystyle\rightarrow$
$\displaystyle x_{0}x_{1}\ldots x_{s-1}$
Clearly from the vertex $X^{s-}$ we can permute all elements in
$\\{x_{s},x_{s+1},\ldots,x_{n-1}\\}$ to determine an out-edge, and at any
point we can permute the $(n-s)$-suffix of a permutation. We would like to
show that, through rotations, the $(n-s)$-suffix may contain any $k$
consecutive blocks from $\\{X_{0},X_{1},\ldots,X_{m-1}\\}$. Consider the
vertex $x_{n-is}x_{n-is+1}\ldots x_{n-(i-1)s-1}$. How does the starting index
$n-is$ correspond to $d$?
$\displaystyle n-is$ $\displaystyle\equiv$ $\displaystyle in-is\pmod{n}$
$\displaystyle\equiv$ $\displaystyle i(n-s)\pmod{n}$ $\displaystyle\equiv$
$\displaystyle ikd\pmod{n}$
Thus the vertices in this rotation cycle always begin with a multiple of $kd$.
Our final step is to show that any multiple of $d$, say $id$, may be written
as a multiple of $kd$ modulo $n$, which is done by Lemma 3.8. Then we can
rotate to start with any block desired, which is equivalent to pushing any $k$
consecutive blocks to the $(n-s)$-suffix.
To summarize, we can perform adjacent transpositions $x_{i}\leftrightarrow
x_{i+1}$ within $X^{s-}$ by rotating $X$ until $x_{i}$ and $x_{i+1}$ fall into
the $k$ blocks in the $(n-s)$-suffix and then transposing. Finally, by
continuing along through rotations we will arrive at the vertex
$x_{0}x_{1}\ldots x_{i-1}x_{i+1}x_{i}x_{i+2}\ldots x_{s-1}$. Thus adjacent
transpositions are always connected, hence all permutations can be reached.
For the converse, suppose that $n-s=d=\hbox{gcd}(n,s)$. In this case,
rotations of the permutation $X=x_{0}x_{1}\ldots x_{n-1}$ provide an
$(n-s)$-suffix of length $d$ – just one block from $X=X_{0}X_{1}\ldots
X_{m-1}$. Thus we can always permute elements within each block, however the
cyclic order of the blocks is fixed and we can perform the swap
$x_{i}\leftrightarrow x_{j}$ if and only if $x_{i}$ and $x_{j}$ are in the
same block. Thus we are able to permute elements within blocks, but not
permute blocks (only rotate the block order). Hence the transition digraph
connects only permutations with the same block order (rotations allowed).
Permutations with block order that are not simple rotations are not connected,
so the graph is disconnected and no Euler tour exists. ∎
Note that Theorem 2.2 agrees with the following well-known fact about
universal cycles, or $(n-1)$-ocycles.
###### Corollary 2.3.
There is no universal cycle for permutations of $M$.
###### Proof.
In this case, we apply Theorem 2.2 with $s=n-1$. Then gcd$(n,s)=1=n-s$, so no
$(n-1)$-ocycle exists. ∎
## 3 Related Objects
### 3.1 Functions
Many times injective, surjective, and bijective functions are represented by
permutations. We have the following facts about functions and theorems
corresponding to their alternate representations.
###### Fact 3.1.
An injective function $f:[k]\rightarrow[n]$ may be represented by
$x_{1}x_{2}\ldots x_{k}$, the $k$-permutation of $[n]$ defined so that
$x_{i}=f(i)$.
###### Theorem 3.2.
_[7]_ For all $n,s,k\in\mathbb{Z}^{+}$ with $1\leq s<k<n$, there is an
$s$-ocycle for $k$-permutations of $[n]$.
###### Fact 3.3.
A surjective function $f:[n]\rightarrow[h]$ may be represented by
$x_{1}x_{2}\ldots x_{n}$, the string with ground set $[h]$ defined so that
$x_{i}=f(i)$.
In [8], it is shown that surjective functions are also represented by weak
orders of $[n]$ with height $h-1$. This observation is followed by the
following theorem.
###### Theorem 3.4.
_[8]_ For all $n,s,h\in\mathbb{Z}^{+}$ with $1\leq s\leq n-2$, gcd$(s,n)=1$,
and $0\leq h\leq n-1$, there is an $s$-ocycle for $\mathcal{W}(n,h)$.
We are able to improve this theorem as follows.
###### Theorem 3.5.
For all $n,s,h\in\mathbb{Z}^{+}$ with $1\leq s\leq n-2$ and $h\leq n-1$ there
is an $s$-ocycle for strings with ground set $[h]$.
###### Proof.
We will show that the corresponding transition graph is eulerian.
Balanced:
Consider a vertex $X^{s-}=x_{1}x_{2}\ldots x_{s}$. $X^{s-}$ is an $s$-prefix
of the string $X=x_{1}x_{2}\ldots x_{n}$ where $X$ has ground set $[h]$. Since
$x_{s+1}x_{s+2}\ldots x_{n}x_{1}x_{2}\ldots x_{s}$ is also a string with
ground set $[h]$, there is a bijection between in- and out-edges at $X^{s-}$.
Hence the graph is balanced.
Connected:
Define the minimum vertex $V^{s-}$ to be the $s$-prefix of the permutation
$V=12\cdots mm\cdots m$. Let $X=x_{1}x_{2}\ldots x_{n}$ be an arbitrary
multiset permutation with ground set $[h]$, and let $X^{s-}$ be the
$s$-prefix. We will show a path from $X^{s-}$ to $V^{s-}$ exists in the
transition graph.
Compare $X^{s-}$ and $V^{s-}$, and let $i$ be the least index such that
$X^{s-}(i)\neq V^{s-}(i)$. Note that since $h\leq n-1$, some element from
$[h]$ must appear at least twice in $X$. We will refer to any element
appearing more than once as a duplicate. We have two cases.
1. 1.
If the letter $x_{i}\in[h]$ appears twice in $X$:
Let $d=\hbox{gcd}(n,s)$, and rotate $X$ until the $d$-block containing $x_{i}$
is first. If we rotate $X$ again by following an out-edge in the graph, we
have arrived at a vertex $A$ representing an $s$-substring of $X$ without
$x_{i}$. Since $x_{i}$ also appears elsewhere in $X$, we can follow
(backwards) the in-edge that is identical to $A$ except with $x_{i}$ replaced
by $v_{i}$. Now we are at an edge corresponding to an $s$-substring of
$x_{1}x_{2}\ldots x_{i-1}v_{i}x_{i+1}\ldots x_{n}$, so we are one step closer
to the minimum vertex.
2. 2.
If the letter $x_{i}\in[h]$ appears exactly once in $X$:
Since $x_{i}$ is not a duplicate in $X$, some other letter $x_{j}\in[h]$ is a
duplicate. In this case, we proceed as in case 1 to replace $x_{j}$ with the
letter $x_{i}$. Then $x_{i}$ is a duplicate so we may follow case 1 again to
replace $x_{i}$ with $v_{i}$. At this point we have moved closer to the
minimum vertex.
Continuing until we have transformed $X^{s-}$ to $V^{s-}$ produces a path from
$X^{s-}$ to the minimum vertex. Hence the graph is connected.
As the transition graph is balanced and connected, it is eulerian by Theorem
1.1. ∎
###### Fact 3.6.
A bijective function $f:[n]\rightarrow[n]$ may be represented by
$x_{1}x_{2}\ldots x_{n}$, the permutation of $[n]$ defined so that
$x_{i}=f(i)$.
###### Theorem 3.7.
Let $n,s\in\mathbb{Z}^{+}$ with $1\leq s\leq n-1$. There exists an $s$-ocycle
on permutations of $[n]$ if and only if $n-s>\hbox{gcd}(n,s)$.
###### Proof.
Use Theorem 2.2 with $M=[n]$. ∎
### 3.2 Juggling Sequences
We begin with some lemmas that will help us to prove our main result.
###### Lemma 3.8.
Let $R=r_{0}r_{1}\ldots r_{n-1}$ be a string, and let $1\leq s\leq n$. Then
$R^{\prime}=\rho^{s}(R)=r_{0-s}r_{1-s}\ldots r_{n-1-s}$ (where addition is
modulo $n$) is a valid juggling sequence if and only if $R$ is a valid
juggling sequence.
###### Proof.
We will check the corresponding permutation sequence for $R^{\prime}$ and show
that it is valid. Suppose for a contradiction that there are
$i,j\in\\{0,1,\ldots,n-1\\}$ with
$r^{\prime}_{i}+i\equiv r^{\prime}_{j}+j\pmod{n}.$
Then we have:
$\displaystyle r^{\prime}_{i}+i$ $\displaystyle\equiv$ $\displaystyle
r^{\prime}_{j}+j\pmod{n}$ $\displaystyle r_{i-s}+i$ $\displaystyle\equiv$
$\displaystyle r_{j-s}+j\pmod{n}$ $\displaystyle r_{i-s}+i-s$
$\displaystyle\equiv$ $\displaystyle r_{j-s}+j-s\pmod{n}$ $\displaystyle
r_{k}+k$ $\displaystyle\equiv$ $\displaystyle r_{\ell}+\ell\pmod{n}$
Thus $\Pi_{R^{\prime}}$ is valid if and only if $\Pi_{R}$ is valid. ∎
###### Lemma 3.9.
Let $T=t_{0}t_{1}\ldots t_{n-1}$ be a juggling sequence, and let $0\leq
i,s\leq n-1$. Then
$\Pi_{\rho^{s}(T)}(i)=\rho^{s}(\Pi_{T})(i)-s\pmod{n}.$
###### Proof.
$\displaystyle\Pi_{\rho^{s}(T)}(i)$ $\displaystyle=$
$\displaystyle\rho^{s}(T)(i)+i\pmod{n}$ $\displaystyle=$ $\displaystyle
T(i+s)+i\pmod{n}$ $\displaystyle=$ $\displaystyle T(i+s)+i+s-s\pmod{n}$
$\displaystyle=$ $\displaystyle\Pi_{T}(i+s)-s\pmod{n}$ $\displaystyle=$
$\displaystyle\rho^{s}(\Pi_{T})(i)-s\pmod{n}$
∎
###### Lemma 3.10.
Fix $n,s,b\in\mathbb{Z}^{+}$ with $1\leq s\leq n-1$. Define $D$ to be the
$s$-ocycle transition digraph for juggling sequences of length $n$ using at
most $b$ balls. From any vertex $v_{0}v_{1}\ldots v_{s-1}$, there exists a
path to any vertex $v_{0}^{\prime}v_{1}^{\prime}\ldots v_{s-1}^{\prime}$ where
$v_{i}^{\prime}\equiv v_{i}\pmod{n}$.
###### Proof.
Since adding/subtracting $n$ to any digit in a juggling sequence of length $n$
does not invalidate the sequence, if $v_{i}\geq n$ we can simply rotate some
juggling sequence $X$ with $s$-prefix $v_{0}v_{1}\ldots v_{s-1}$ until $v_{i}$
is in the $(n-s$)-suffix, replace with $v_{i}$ with $v_{i}-n$, and then rotate
back to our original $s$-prefix with $v_{i}$ replaced by $v_{i}-n$. Repeating
this eventually will find a path to $v_{1}^{\prime}v_{2}^{\prime}\ldots
v_{s}^{\prime}$. ∎
###### Theorem 3.11.
Fix $n,s,b\in\mathbb{Z}^{+}$ such that $1\leq s\leq n-1$. There exists an
$s$-ocycle for the set of juggling sequences with period $n$ and at most $b$
balls if and only if $n-s>\hbox{gcd}(n,s)$.
###### Proof.
We prove the forward direction by showing that the $s$-ocycle transition
digraph $D$ has an Euler tour. By Theorem 1.1 this is done by showing that the
graph is balanced and connected.
Balanced:
Consider a vertex $R=r_{0}r_{1}\ldots r_{s-1}$ in $D$. We will show that any
$(n-s)$-string $Q=q_{0}q_{1}\ldots q_{n-s-1}$ that is a valid $s$-suffix for
$R$ is also a valid $s$-prefix for $R$. Note that strings $RQ$ and $QR$ are
simply rotations of each other, so by Lemma 3.8 either both strings are valid
juggling sequences or neither string is. In this manner, there is a bijection
between in- and out-edges, hence all vertices are balanced.
Connected:
Consider an arbitrary vertex $T^{s-}=t_{0}t_{1}\ldots t_{s-1}$ that is an
$s$-prefix to some juggling sequence $T=t_{0}t_{1}\ldots t_{n-1}$. First, by
Lemma 3.10 we may assume that $t_{i}\in\\{0,1,\ldots,n-1\\}$ for all
$i\in\\{0,1,\ldots,n-1\\}$. We will show that this arbitrary vertex is
connected to the minimum vertex, which we define to be $V^{s-}=0^{s}$ (a
prefix of $V=0^{n}$). In doing so, we will have shown that every vertex is
connected to $V^{s-}$ and hence the graph is connected.
Compare permutation sequences $\Pi_{T}$ and $\Pi_{V}$ corresponding to
juggling sequences $T$ and $V$, respectively. Note that
$\Pi_{V}=012\ldots(n-1)$, so suppose that for all $x\in\\{0,1,\ldots,i-1\\}$
we have $\Pi_{T}(x)=\Pi_{V}(x)=x$, but that $\Pi_{T}(i)=j$ for some
$j\in\\{i+1,i+2,\ldots,s-1\\}$. We will find a path from the $s$-prefix of $T$
to the $s$-prefix of some juggling sequence $T^{\prime}$ with permutation
sequence $\Pi_{T^{\prime}}$ that agrees with $\Pi_{V}$ in the first $i+1$
positions. Repeating this procedure until $i=n$ will construct a path through
$D$ from $T^{s-}$ to $V^{s-}$.
Assume $n-s=kd$ and $n=md$ for integers $m,k$ where $d=\hbox{gcd}(n,s)$, and
let $T=Y_{0}Y_{1}\ldots Y_{m-1}$ be a partition of $T$ into $d$-blocks, i.e.
$Y_{a}=t_{ad}t_{ad+1}\ldots t_{ad+d-1}$. Suppose that $t_{i}\in Y_{a}$ and
$t_{j}\in Y_{b}$. We have two cases.
1. 1.
If $b-a<k$:
In this case, from Lemma 3.8 we may perform $s$-rotations on $T$ so that we
arrive at a vertex with both $t_{i}$ and $t_{j}$ in the $(n-s)$-suffix of $T$.
Suppose that we performed a total rotation of size $R$, i.e. we are now at the
vertex that represents the $s$-prefix of $\rho^{R}(T)$.
At this point, the values $t_{i}$ and $t_{j}$ are located in positions $i-R$
and $j-R$ of $\rho^{R}(T)$, respectively. By Lemma 3.9 the corresponding
permutation sequence entries are:
$\displaystyle\Pi_{\rho^{R}(T)}(i-R)$ $\displaystyle=$
$\displaystyle\rho^{R}(\Pi_{T})(i-R)-R\pmod{n}$ $\displaystyle=$
$\displaystyle\Pi_{T}(i)-R\pmod{n}$
and
$\displaystyle\Pi_{\rho^{R}(T)}(j-R)$ $\displaystyle=$
$\displaystyle\Pi_{T}(j)-R\pmod{n}.$
We now note that the vertex $(\rho^{R}(T))^{s-}$, defined as the $s$-prefix of
$\rho^{R}(T)$, is also the $s$-prefix of the juggling sequence
$\rho^{R}(T^{\prime})$ that is obtained from $\rho^{R}(T)$ by performing the
swap $\Pi_{\rho^{R}(T)}(i-R)\leftrightarrow\Pi_{\rho^{R}(T)}(j-R)$ and
adjusting the values $\rho^{R}(T)(i-R)$ and $\rho^{R}(T)(j-R)$ appropriately.
Then rotating backwards by $R$ we reach a juggling sequence $T^{\prime}$ with
$\Pi_{T^{\prime}}(x)=x$ for all $x\in\\{0,1,\ldots,i-1,i\\}$. This means that
we have found a path from $T^{s-}$ to $T^{\prime s-}$, which is one step
closer to the minimum vertex.
2. 2.
If $b-a\geq k$:
In this case the blocks containing $t_{i}$ and $t_{j}$ are more than $k$
apart, so we cannot rotate to have both $t_{i}$ and $t_{j}$ in the
$(n-s)$-suffix of $T$ simultaneously. Instead, we pick a point $t_{z}$ in the
block $Y_{c}$ that is $k$ blocks preceding $Y_{b}$ and follow case (1) with
$t_{z}$ in place of $t_{i}$. This produces a juggling sequence $T^{\prime}$
with $\Pi_{T^{\prime}}(z)=i$, where $i\leq z<j$. Repeating, we will eventually
have moved the permutation sequence value $i$ to a block that is close enough
to block $Y_{a}$ to apply Case (1).
By repeating the above procedures, we will eventually have transitioned to the
vertex corresponding to a juggling sequence with $s$-prefix $01\ldots(s-1)$.
At this point we have reached the minimum vertex and we are done.
For the reverse direction, suppose that $n-s=\hbox{gcd}(n,s)=d$. Recall that
the number of balls $b\in\mathbb{Z}^{+}$ is given by:
$b=\frac{1}{n}\sum_{i=0}^{n-1}t_{i}.$
Thus juggling sequences of period $n$ must always satisfy
$\sum_{i=0}^{n-1}t_{i}\equiv 0\pmod{n}$. Equivalently, when we partition a
juggling sequence $T=t_{0}t_{1}\ldots t_{n-1}$ into blocks of length $d$, i.e.
$T=Y_{0}Y_{1}\ldots Y_{m-1}$, where $n=md$, we must have:
$\sum_{i=0}^{m-1}w(Y_{i})=\sum_{i=0}^{m-1}\sum_{j=0}^{d}t_{id+j}=\sum_{i=0}^{n-1}t_{i}\equiv
0\pmod{n},$
where, for a given block $Y_{i}$, we call $w(Y_{i})=\sum_{j=0}^{d}t_{id+j}$
the weight of block $Y_{i}$.
Now since $n-s=\hbox{gcd}(n,s)$, from vertex $T^{s-}=t_{0}t_{1}\ldots t_{s-1}$
in $D$ we may only move to vertices in which the $(n-s)$-suffix has weight
equivalent to $w(Y_{m-1})$ mod $n$. Thus if we can show that, for any $n,s,$
and $b$, there exists a juggling sequence with a block with weight
$w\not\equiv 0\pmod{n}$, then we are done. This is witnessed by the juggling
sequence
$T=\begin{array}[]{ccccccccc}d&0&0&\ldots&0&(n-d)&0&\ldots&0,\end{array}$
with permutation sequence
$\Pi_{T}=\begin{array}[]{cccccccccc}d&1&2&\ldots&(d-1)&0&(d+1)&(d+2)&\ldots&(n-1).\end{array}$
This juggling sequence utilizes one ball (recall we required
$b\in\mathbb{Z}^{+}$) and the weight of the first block of length $d$ is
$d\not\equiv 0\pmod{n}$. Thus this first block must always have weight equal
to $d$ modulo $n$ if $n-s=d$, and so the vertex $T^{s-}$ that represents the
$s$-prefix of $T$ is not connected to the vertex $0^{s}$. Hence no Euler tour
can exist and so no $s$-ocycle exists. ∎
While Theorem 3.11 completes the question of when $s$-ocycles for juggling
sequences of period $n$ and at most $b$ balls exist, several variations remain
open. For example, one might consider juggling sequences with:
* •
exactly $b$ balls,
* •
at least $b$ balls,
* •
fixed minimum period, or
* •
period $n$ with no restriction on number of balls.
## References
* [1] J. Buhler, D. Eisenbud, R. Graham, and C. Wright, _Juggling Drops and Descents_ , Amer. Math. Monthly 101:6, 507-519.
* [2] F. Chung, P. Diaconis, and R. Graham, _Universal Cycles for Combinatorial Structures_ , Discrete Mathematics 110 (1992), 43-59.
* [3] F. Chung and R. Graham, _Universal Juggling Cycles_ , Integers: Electronic Journal of Combinatorial Number Theory 7:2 (2007), #A08.
* [4] S. Forrest, _Genetic Algorithms: Principles of Natural Selection Applied to Computation_ , Science 261:5123 (1993), 872-878.
* [5] A. Godbole, D. Knisley, and R. Norwood, _On $\alpha$-Overlap Graphs_, Proceedings of the Forty-First Southeastern International Conference on Combinatorics, Graph Theory and Computing, Congr. Numer. 204 (2010), 161-171.
* [6] F. Gray, _Pulse code communication_ , March 17, 1953 (filed Nov. 1947), U.S. Patent 2,632,058.
* [7] V. Horan and G. Hurlbert, _$s$ -Overlap Cycles for Permutations,_ to appear in Bull. Inst. Combin. Appl.
* [8] V. Horan and G. Hurlbert, _Universal Cycles for Weak Orders,_ To appear in SIAM Journal of Discrete Math.
* [9] G. Hurlbert and G. Isaak, _Equivalence class universal cycles for permutations_ , Discrete Math. 149 (1996), pp. 123–129.
* [10] J. R. Johnson, _Universal cycles for permutations_ , Discrete Math. 309 (2009), pp. 5264–5270.
* [11] E. A. Ragland, V. Nuys, and H. B. Scultheis, Jr., _Direction-Sensitive Binary Code Position Control System_ , Feb. 11, 1958 (filed Oct. 1953), U.S. Patent 2,823,345.
* [12] Z. Wang and J. Bruck, _Partial Rank Modulation for Flash Memories_ , IEEE International Symposium on Information Theory Proceedings (2010), 864-868.
* [13] D. B. West, _Introduction to Graph Theory_. Prentice Hall, Second Edition, 2001\.
|
arxiv-papers
| 2013-09-18T18:39:30 |
2024-09-04T02:49:51.130395
|
{
"license": "Public Domain",
"authors": "Victoria Horan",
"submitter": "Victoria Horan",
"url": "https://arxiv.org/abs/1309.4741"
}
|
1309.4953
|
# A Near Optimal Approximation Algorithm for Vertex-Cover Problem
Deepak [email protected]
Department of Computer Science & Engineering
National Institute of Technology Silchar
Silchar, Assam, India
###### Abstract
Recently, there has been increasing interest and progress in improvising the
approximation algorithm for well-known NP-Complete problems, particularly the
approximation algorithm for the Vertex-Cover problem. Here we have proposed a
polynomial time efficient algorithm for vertex-cover problem for more
approximate to the optimal solution, which lead to the worst time complexity
$\Theta(V^{2})$ and space complexity $\Theta(V+E)$. We show that our proposed
method is more approximate with example and theorem proof. Our algorithm also
induces improvement on previous algorithms for the independent set problem on
graphs of small and high degree.
Keywords:Approximation algorithm Vertex-Cover Problem Complexity Adjacency
list.
## 1 Introduction
A graph G represents as G = (V, E): V is number of vertices in graph and E is
number of edges in graph and impliment as an adjacency lists or as an
adjacency matrix for both directed and undirected graphs. There are two types
of graph i.e. (I) Sparse graphs-those for which $|E|$ is much less than
$|V|^{2}$ $(E<<V^{2})$. (II) Dense graphs-those for which $|E|$ is close to
$|V|^{2}(E\simeq V^{2})$. Here we presented the graph as the adjacency-list
for the evaluation of our algorithm.
The vertex-cover problem is to find a minimum number of vertex to cover a
given undirected graph. We call such a vertex cover an optimal vertex cover.
This problem is the optimization version of an NP-complete decision problem.
Proposed algorithm is polynomial time algorithm in order to find the set of
vertex to cover the graph. Which shows the better performance than the
traditional algorithm for vertex cover [1].
## 2 The vertex-cover problem
As it is NP-Hard problem so it is hard to find an optimal solution of a graph
G, but not difficult to find a near optimal solution. Our propose method gives
very near optimal solution for Vertex-cover problem. The following
approximation algorithm takes an undirected graph G as input [1] and returns a
set of vertex to cover the graph and whose size is less than the previous
method.
All graphs mentioned here are simple undirected graph. We follow [2] for
definitions. Our proposed method (See algorithm and Fig.2) is on undirected
graph. Here we used the adjacency list to represent graph G. We introduce a
new field weight in the list to store the degree of each individual vertex.
i.e.
struct list
{
char vertex;
int weight;
struct node *next;
struct node *ref;
};
Algorithm 1 Approximate Vertex-Cover Algorithm
0: In the List we introduce another field weight The value of weight is number
of node in reference (ref)
1: $C^{+}\leftarrow\emptyset$
2: L = List
3: L[w] = Reference weight
4: (h, v) = highest weight of the list and respective vertex
5: if $h\neq 0$ then
6: $C^{+}\leftarrow C^{+}\cup{v}$
7: v[w] $\leftarrow 0$
8: for all vertex of List L[ref] $\in\\{v\\}$ do
9: L[w] $\leftarrow$ L[w]-1
10: end for
11: go to 4
12: else
13: return $C^{+}$
14: end if
It’s space complexity is $\Theta(V+E)$ [1]. For step 4 search in the list is
O(V). In the for loop i.e. step 8 to 10 for each individual vertex need to
search its reference vertices. So worst time complexity is $O(V*(V-1))$=
$O(V^{2})$. So the worst time complexity of the graph is $\Theta(V^{2})$.
###### Theorem 1.
_(Thomas H. Cormen et.al.[1])_ APPROX-VERTEX-COVER is a polynomial-time
2-approximation algorithm.
###### Theorem 2.
Proposed Approximate Vertex-Cover is a polynomial-time
$(2-\varepsilon)$-approximation algorithm.
###### Proof.
In Theorem 1 $C$ is the set of vertex for APPROX-VERTEX-COVER and $C^{*}$ in
the optimal vertex cover $i.e.|C|\leq 2|C^{*}|$. In our approach we pick one
vertex and remove the edges connected to that vertex. So most of the times we
don’t consider both end point of one edges, which followed in Theorem 1. For
our proposed method we consider the resultant set of vertex is $C^{+}$, then
$|C^{+}|=|C|-\varepsilon$.
$\Rightarrow|C^{+}|=(2-\varepsilon)|C^{*}|$, $0\leq\varepsilon\leq 1$. ∎∎
In some cases proposed method$(C^{+})$ approaches to optimal solution when
$\varepsilon$ value is 1. Our method is shown in Fig. 1, algorithm
(Approximate Vertex-Cover), and proved in Theorem 2. The comparison of the
optimal vertex-cover, previous vertex-cover and proposed vertex-cover shown in
Fig. 2.
Figure 1: The operation of Approximate Vertex-Cover. (a) The input graph G,
which has 7 vertices and 8 edges. (b) The vertex D has the highest weight, is
the first vertex chosen by Approximate Vertex-Cover. The weight of the vertex
initiate to zero and the vertices directly connected from vertex D weight
value decrees by 1. The vertex D is shaded and the edges is dashed incident
from vertex D. The vertex D add to the set $C^{+}$. (c) There are three vertex
with highest weight (B, C, E); Arbitrary vertex C is chosen and add to the set
$C^{+}$. (d) Vertex A is chosen; added to $C^{+}$. (e) Vertex F is chosen;
added to $C^{+}$. (f) The resultant vertex cover with our proposed method.
Approximate Vertex-Cover, contains the four vertices a, c, d, f. Figure 2:
Comparison of optimal, Traditional and proposed vertex cover result. (a) The
inputted original graph G. (b) The optimal result of input graph. (c) The
result with traditional method. (d) The proposed algorithm’s result for
vertex-cover algorithm.
## 3 Conclusion
Here in our proposed technique we produce the set of vertex for vertex-cover
problem. Which is more near optimal solution and better than the previous
technique.
## References
* [1] Thomas H. Cormen et.al. Introduction to Algorithms, Second Edition. The MIT Press, 2001. ISBN 0-262-03293-7
* [2] Liu, Chung Laung. Elements of discrete mathematics. Vol. 101. New York: McGraw-Hill, 1985.
|
arxiv-papers
| 2013-09-19T12:27:27 |
2024-09-04T02:49:51.151682
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Deepak Puthal",
"submitter": "Deepak Puthal",
"url": "https://arxiv.org/abs/1309.4953"
}
|
1309.4986
|
# Four generated, squarefree, monomial ideals
Adrian Popescu Adrian Popescu, Department of Mathematics, University of
Kaiserslautern, Erwin-Schrödinger-Str., 67663 Kaiserslautern, Germany
[email protected] and Dorin Popescu Dorin Popescu, Simion Stoilow
Institute of Mathematics of Romanian Academy, Research unit 5, P.O.Box 1-764,
Bucharest 014700, Romania [email protected]
###### Abstract.
$I\supsetneq J$ be two squarefree monomial ideals of a polynomial algebra over
a field generated in degree $\geq d$, resp. $\geq d+1$ . Suppose that $I$ is
either generated by three monomials of degrees $d$ and a set of monomials of
degrees $\geq d+1$, or by four special monomials of degrees $d$. If the
Stanley depth of $I/J$ is $\leq d+1$ then the usual depth of $I/J$ is $\leq
d+1$ too.
Monomial Ideals, Depth, Stanley depth.
2010 Mathematics Subject Classification: Primary 13C15, Secondary 13F20,
13F55, 13P10.
The support of the first author from the Department of Mathematics of the
University of Kaiserslautern and the support of the second author from grant
PN-II-RU-TE-2012-3-0161 of Romanian Ministry of Education, Research and
Innovation are gratefully acknowledged.
## Introduction
Let $K$ be a field and $S=K[x_{1},\ldots,x_{n}]$ be the polynomial $K$-algebra
in $n$ variables. Let $I\supsetneq J$ be two squarefree monomial ideals of $S$
and suppose that $I$ is generated by squarefree monomials of degrees $\geq d$
for some positive integer $d$. After a multigraded isomorphism we may assume
either that $J=0$, or $J$ is generated in degrees $\geq d+1$. By [5,
Proposition 3.1] (see [12, Lemma 1.1]) we have
$\operatorname{depth}_{S}I/J\geq d$. Depth of $I/J$ is a homological invariant
and depends on the characteristic of the field $K$.
The purpose of our paper is to study upper bound conditions for
$\operatorname{depth}_{S}I/J$. Let $B$ (resp. $C$) be the set of the
squarefree monomials of degrees $d+1$ (resp. $d+2$) of $I\setminus J$. Suppose
that $I$ is generated by some squarefree monomials $f_{1},\ldots,f_{r}$ of
degrees $d$ for some $d\in{\mathbb{N}}$ and a set of squarefree monomials $E$
of degree $\geq d+1$. If $d=1$ and each monomial of $B\setminus E$ is the
least common multiple of two $f_{i}$ then it is easy to show that
$\operatorname{depth}_{S}I/J=1$ (see Lemma 3). Trying to extend this result
for $d>1$ we find an obstruction given by Example 2. Our extension given by
Lemma 4 is just a special form, but a natural condition seems to be given in
terms of the Stanley depth.
More precisely, let $P_{I\setminus J}$ be the poset of all squarefree
monomials of $I\setminus J$ with the order given by the divisibility. Let $P$
be a partition of $P_{I\setminus J}$ in intervals $[u,v]=\\{w\in P_{I\setminus
J}:u|w,w|v\\}$, let us say $P_{I\setminus J}=\cup_{i}[u_{i},v_{i}]$, the union
being disjoint. Define
$\operatorname{sdepth}P=\operatorname{min}_{i}\operatorname{deg}v_{i}$ and the
Stanley depth of $I/J$ given by
$\operatorname{sdepth}_{S}I/J=\operatorname{max}_{P}\operatorname{sdepth}P$,
where $P$ runs in the set of all partitions of $P_{I\setminus J}$ (see [5],
[20]). Stanley’s Conjecture says that
$\operatorname{sdepth}_{S}I/J\geq\operatorname{depth}_{S}I/J$. The Stanley
depth of $I/J$ is a combinatorial invariant and does not depend on the
characteristic of the field $K$. Stanley’s Conjecture holds when $J=0$ and $I$
is an intersection of four monomial prime ideals by [8], [10], or $I$ is such
that the sum of every three different of its minimal prime ideals is a
constant ideal by [11] (see also [14]), or $I$ is an intersection of three
monomial primary ideals by [22], or a monomial almost complete intersection by
[4].
###### Theorem 1.
(D. Popescu [12, Theorem 4.3]) If $\operatorname{sdepth}_{S}I/J=d$ then
$\operatorname{depth}_{S}I/J=d$, that is Stanley’s Conjecture holds in this
case.
Next step in the study of Stanley’s Conjecture is to show the following weaker
conjecture.
###### Conjecture 1.
Suppose that $I\subset S$ is minimally generated by some squarefree monomials
$f_{1},\ldots,f_{r}$ of degrees $d$, and a set $E$ of squarefree monomials of
degrees $\geq d+1$. If $\operatorname{sdepth}_{S}I/J=d+1$ then
$\operatorname{depth}_{S}I/J\leq d+1$.
Set $s=|B|$, $q=|C|$. In the study of the above conjecture very useful seem to
be the following two particular results of [13, Theorem 1.3] and [19, Theorem
2.4].
###### Theorem 2.
(D. Popescu) If $s>q+r$ then $\operatorname{depth}_{S}I/J\leq d+1$.
###### Theorem 3.
(Y. Shen) If $s<2r$ then $\operatorname{depth}_{S}I/J\leq d+1$.
These results were hinted by Stanley’s Conjecture since it is obvious that
$s>q+r$, or $s<2r$ imply $\operatorname{sdepth}_{S}I/J\leq d+1$. The proof of
Theorem 2 uses Koszul homology (see [1, Section 1.6]). Shen’s proof of the
above theorem as well of Theorem 2 is easy and uses the Hilbert depth
considered by Bruns-Krattenhaler-Uliczka [2] (see also [21], [6]).
An equivalent definition for the Stanley depth is:
$\operatorname{sdepth}(M)=\operatorname{max}\\{\operatorname{sdepth}\mathcal{D}\
|\ \mathcal{D}\textnormal{ is a Stanley decomposition of }M\\},$
where a Stanley decomposition of a $\mathbb{Z}-$graded (resp.
$\mathbb{Z}^{n}-$graded) $S-module$ $M$ is $\mathcal{D}=(S_{i},u_{i})_{i\in
I}$, where $u_{i}$ are homogenous elements of $M$ and $S_{i}$ are graded
(resp. $\mathbb{Z}^{n}-$graded) $K-$algebra retracts of $S$ and
$S_{i}\cap\operatorname{Ann}(u_{i})=0$ such that $M=\oplus_{i}S_{i}u_{i}$; and
$\operatorname{sdepth}\mathcal{D}$ is the $\operatorname{depth}$ of the
$S-$module $\oplus_{i}S_{i}u_{i}$. A more general concept is the one of
Hilbert depth of a $\mathbb{Z}-$graded module $M$, denoted by
$\operatorname{hdepth}_{1}(M)$. Instead of considering equality, we only
assume that $M\cong\oplus S_{i}(-s_{i})$, where $s_{i}\in\mathbb{Z}$. One can
also construct $\operatorname{hdepth}_{n}$ analogously if $M$ is a multigraded
(that is $\mathbb{Z}^{n})$ module.
In [9] is presented (and implemented) an algorithm that computes
$\operatorname{hdepth}_{1}(M)$ based on a Theorem of Uliczka [21]; and in [7]
was presented an algorithm that computes $\operatorname{hdepth}_{n}(M)$.
Meanwhile, another algorithm that computes $\operatorname{hdepth}_{1}$ and
more was given in [3]. [9, Proposition 1.9] gives a partial answer to a
question of Herzog asking whether
$\operatorname{sdepth}m=\operatorname{sdepth}(S\oplus m)$, where $m$ is the
graded maximal ideal of $S$. More precisely, for $n\in\\{1,2,3,4,5,7,9,11\\}$
one obtains $\operatorname{hdepth}_{1}m=\operatorname{hdepth}_{1}(S\oplus m)$,
which gives $\operatorname{sdepth}m=\operatorname{sdepth}(S\oplus m)$ (again
Hilbert depth helps the study of Stanley depth). For $n=6$ we have
$\operatorname{hdepth}_{1}m\neq\operatorname{hdepth}_{1}(S\oplus m)$, which
means that in general Herzog’s question could have a negative answer. Later
Ichim and Zarojanu checked the case $n=6$ and found indeed a counterexample to
Herzog’s question, which will be included in the new version of [7].
An important step in proving Conjecture 1 is the following theorem.
###### Theorem 4.
(D. Popescu-A. Zarojanu [16], [17, Theorem 1.5]) Conjecture 1 holds in each of
the following two cases:
1. (1)
$r=1$,
2. (2)
$1<r\leq 3$, $E=\emptyset$.
Next theorem is the main result of this paper.
###### Theorem 5.
Conjecture 1 holds in each of the following two cases:
1. (1)
$r\leq 3$,
2. (2)
$r=4$, $E=\emptyset$ and there exists $c\in C$ such that
$\operatorname{supp}c\not\subset\cup_{i\in[4]}\operatorname{supp}f_{i}$.
This follows from our Theorems 6, 8. The proof of 6 extends the proof of [17,
Theorem 2.3].
We owe thanks to A. Zarojanu, who noticed some small mistakes in a previous
version of this paper and gave us the bad example 5.
## 1\. Depth and Stanley depth
Let $I\supsetneq J$ be two squarefree monomial ideals of $S$. We assume that
$I$ is generated by squarefree monomials $f_{1},\ldots,f_{r}$ of degrees $d$
for some $d\in{\mathbb{N}}$ and a set of squarefree monomials $E$ of degree
$\geq d+1$. We may suppose that either $J=0$, or is generated by some
squarefree monomials of degrees $\geq d+1$. $B$ (resp. $C$) denotes the set of
the squarefree monomials of degrees $d+1$ (resp. $d+2$) of $I\setminus J$.
###### Lemma 1.
Let $J\subset I$ be square free monomial ideals and $j\in[n]$ be such that
$(J:x_{j})\not=(I:x_{j})$. Then
$\operatorname{depth}_{S}(I:x_{j})/(J:x_{j})\geq\operatorname{depth}_{S}I/J$.
###### Proof.
We have
$\operatorname{pd}_{S}I/J\geq\operatorname{pd}_{S_{x_{j}}}(I/J)\otimes
S_{x_{j}}=\operatorname{pd}_{S_{x_{j}}}((I:x_{j})/(J:x_{j}))\otimes
S_{x_{j}}=\operatorname{pd}_{S}((I:x_{j})/(J:x_{j}))$
the last equality holds since $x_{j}$ does not appear among the generators of
$(I:x_{j})$ and $(J:x_{j})$. Now it is enough to apply the Auslander-Buchsbaum
Theorem.
###### Lemma 2.
Let $t\in[n]$. Suppose that $I\not=J+I\cap(x_{t})$ and
$\operatorname{depth}_{S}I/(J+I\cap(x_{t}))=d$. If
$\operatorname{depth}_{S}I/J\geq d+1$ then $\operatorname{depth}_{S}I/J=d+1$.
###### Proof.
In the following exact sequence
$0\rightarrow(I:x_{t})/(J:x_{t})\xrightarrow{x_{t}}I/J\rightarrow
I/(J+I\cap(x_{t}))\rightarrow 0$
the first term has depth $d+1$ by the Depth Lemma. Now it is enough to apply
the above lemma.
Let $w_{ij}$ be the least common multiple of $f_{i}$ and $f_{j}$ and set $W$
to be the set of all $w_{ij}\in B$.
###### Lemma 3.
If $d=1$ and $B\subset E\cup W$ then $\operatorname{depth}_{S}I/J=1$.
###### Proof.
First suppose that $E=\emptyset$, let us say $I=(x_{1},\ldots,x_{r})$. Set
$S^{\prime}=K[x_{1},\ldots,x_{r}]$, $I^{\prime}=I\cap S^{\prime}$,
$J^{\prime}=J\cap S^{\prime}$. By hypothesis $B\subset S^{\prime}$ and it
follows that $(x_{r+1},\ldots,x_{n})I\subset J$ and so
$\operatorname{depth}_{S}I=\operatorname{depth}_{S^{\prime}}I^{\prime}=1$. But
$\operatorname{depth}_{S}J\geq 2$, if $J\not=0$, and so
$\operatorname{depth}_{S}I/J=1$ by the Depth Lemma.
Now, suppose that $E\not=\emptyset$. In the following exact sequence
$0\rightarrow(x_{1},\ldots,x_{r})/J\cap(x_{1},\ldots,x_{r})\rightarrow
I/J\rightarrow I/(J,x_{1},\ldots,x_{r})\rightarrow 0$
the first term has depth 1 as above and the last term has depth $\geq d+1$
since it is generated by squarefree monomials of degrees $\geq 2$ from $E$.
Again the Depth Lemma gives $\operatorname{depth}_{S}I/J=1$.
###### Lemma 4.
Suppose that $I\subset S$ is generated by some squarefree monomials
$f_{1},...f_{r}$ of degree $d$. Assume that for all $b\in B$ all divisors of
$b$ of degree $d$ are among $\\{f_{1},\ldots,f_{r}\\}$. Then
$\operatorname{depth}_{S}I/J=d$.
###### Proof.
Apply induction on $d\geq 1$. If $d=1$ then apply the above lemma. Assume
$d>1$. We may suppose that $n\in\operatorname{supp}f_{1}$. $(I:x_{n})$ is an
extension of a squarefree monomial ideal $I^{\prime}$ of
$S^{\prime}=K[x_{1},\ldots,x_{n-1}]$ which is generated in degree $\geq d-1$.
Similarly $(J:x_{n})$ is generated by a squarefree monomial ideal $J^{\prime}$
of $S^{\prime}$. Note that the generators of $I^{\prime}$ of degree $d-1$ have
the form $f^{\prime}_{i}=f_{i}/x_{n}$ for $f_{i}\in(x_{n})$, and the
squarefree monomials $B^{\prime}$ of degrees $d$ from $I^{\prime}\setminus
J^{\prime}$ have the form $b^{\prime}=b/x_{n}$ for some $b\in(B\cap(x_{n}))$.
Certainly we must consider also the case when $f_{j}\not\in(x_{n})$. If
$x_{n}f_{j}\in J$ then $f_{j}\in(J:x_{n})$ is not in $B^{\prime}$. Otherwise,
$f_{j}=(x_{n}f_{j})/x_{n}\in B^{\prime}$. Note that all divisors of degree
$d-1$ of each $b^{\prime}\in B^{\prime}$ are among $f^{\prime}_{i}$. By
induction hypothesis we have
$\operatorname{depth}_{S^{\prime}}I^{\prime}/J^{\prime}=d-1$ and so
$\operatorname{depth}_{S}(I:x_{n})/(J:x_{n})=d$. Now it is enough to apply
Lemma 1.
An obstruction to improve Lemma 3 and the above lemma is given by the
following example.
###### Example 1.
Let $n=5$, $d=2$, $r=5$,
$I=(x_{1}x_{2},x_{1}x_{3},x_{2}x_{3},x_{1}x_{4},x_{3}x_{5})$,
$J=(x_{1}x_{2}x_{5},x_{1}x_{4}x_{5},x_{2}x_{3}x_{4},x_{3}x_{4}x_{5})$,
$B=\\{x_{1}x_{2}x_{3},x_{1}x_{2}x_{4},x_{1}x_{3}x_{4},x_{1}x_{3}x_{5},x_{2}x_{3}x_{5}\\}$.
We have $\operatorname{depth}_{S}I/J=3$ because
$\operatorname{depth}_{S}S/J=3$, $\operatorname{depth}_{S}S/I=2$ and with the
help of Depth Lemma. Note that each $b\in B$ is the least common multiple of
two generators of $I$, but for example $b=x_{1}x_{2}x_{4}$ has
$x_{2}x_{4}\not\in I$ as a divisor of degree $2$.
Let $C_{2}=C\cap W$ and $C_{3}$ be the set of all $c\in C$ having all divisors
from $B\setminus E$ in $W$. In particular each monomial of $C_{3}$ is the
least common multiple of three of $f_{i}$. The converse is not true as shows
the following example.
###### Example 2.
Let $n=4$, $d=2$, $r=3$, $f_{1}=x_{1}x_{2}$, $f_{2}=x_{2}x_{3}$,
$f_{3}=x_{3}x_{4}$, $I=(f_{1},f_{2},f_{3})$ and $J=0$. Then
$c=x_{1}x_{2}x_{3}x_{4}$ is the least common multiple of $f_{1},f_{2},f_{3}$
but has a divisor $b=x_{1}x_{2}x_{4}\in B$ which is not the least common
multiple of two $f_{i}$.
Next theorem is our key result, its proof is based on [17, Theorem 2.1] and
will be given in the last section. The main reason that this proof works for
$r\leq 3$ but not for $r=4$ is that in the first case $|C_{3}|\leq 1$ but in
the second one we may have $|C_{3}|=4$, which makes the things harder.
However, for $r\geq 5$ will appear a new problem since we may have $B\subset
W$ and $s\geq 2r$ (for example when $r=5$, $d=2$ we may have $s=10=2r$). We
remind that by Theorem 3 we had to check Stanley’s Conjecture only when $s\geq
2r$.
###### Theorem 6.
Conjecture 1 holds for $r\leq 3$, the case $r=1$ being given in Theorem 4.
###### Example 3.
Let $n=5$, $f_{1}=x_{1}x_{2}$, $f_{2}=x_{1}x_{3}$, $f_{3}=x_{1}x_{4}$,
$a=x_{2}x_{3}x_{5}$, $E=\\{a\\}$, $I=(f_{1},f_{2},f_{3},a)$, $J=(x_{4}a)$. We
have $w_{12}=f_{1}x_{3}$, $w_{13}=f_{1}x_{4}$, $w_{23}=f_{2}x_{4}$. Set
$c=w_{12}x_{4}$, $c_{1}=w_{12}x_{5}$, $c_{2}=w_{23}x_{5}$,
$c_{3}=w_{13}x_{5}$. Then $C=\\{c,c_{1},c_{2},c_{3}\\}$ and $B\setminus
E=B\cap(\cup_{i}[f_{i},c_{i}])$. Thus $s=7$, $q=4$, $r=3$. It is easy to see
that $\operatorname{sdepth}_{S}I/J=3$. Indeed, note that $c_{1}$ is the only
$c^{\prime}\in C$ which is multiple of $a$. Suppose that there exists a
partition $P$ on $P_{I/J}$ with sdepth $4$. Then we have necessarily in $P$
the interval $[a,c_{1}]$. If $P$ contains the interval $[f_{1},c]$ then it
must contain also the intervals $[f_{2},c_{2}]$ and so $[f_{3},c_{3}]$, but
then $w_{13}\in[f_{1},c]\cap[f_{3},c_{3}]$, that is the union is not disjoint.
If $P$ contains the interval $[f_{1},c_{3}]$ then $P$ contains either
$[f_{3},c]$, $[f_{2},c_{2}]$, or $[f_{2},c]$, $[f_{3},c_{2}]$, in both cases
the intersection of these two intervals contains $w_{23}$, which is false. By
Theorem 6 we get $\operatorname{depth}_{S}I/J\leq 3$, this inequality being in
fact an equality.
## 2\. A special case of $r=4$
###### Theorem 7.
Suppose that $I\subset S$ is minimally generated by some squarefree monomials
$\\{f_{1},\ldots,f_{r}\\}$ of degrees $d$ such that there exists $c\in C$ with
$\operatorname{supp}c\not\subset\cup_{i\in[r]}\operatorname{supp}f_{i}$. If
Conjecture 1 holds for $r^{\prime}<r$ and $\operatorname{sdepth}_{S}I/J=d+1$,
then $\operatorname{depth}_{S}I/J\leq d+1$.
###### Proof.
By [17, Lemma 1.1] we may assume that $C\subset(W)$. By hypothesis, choose
$t\in\operatorname{supp}c$ such that
$t\not\in\cup_{i\in[r]}\operatorname{supp}f_{i}$. We may suppose that
$B\cap(x_{t})=\\{x_{t}f_{1},\ldots x_{t}f_{e}\\}$ for some $1\leq e\leq r$.
Set $I_{t}=I\cap(x_{t})$, $J_{t}=J\cap(x_{t})$ and $U_{t}=I_{t}/J_{t}$. Then
$B_{t}$ generates $I_{t}$.
First assume that $\operatorname{sdepth}_{S}U_{t}\leq d+1$. It follows that
$\operatorname{depth}_{S}U_{t}\leq d+1$ by [12, Theorem 4.3]. But
$U_{t}\cong(I:x_{t})/(J:x_{t})$ and so
$\operatorname{depth}_{S}U_{t}\geq\operatorname{depth}_{S}I/J$ by Lemma 1,
which is enough.
Now assume that $U_{t}$ has sdepth $\geq d+2$. Let $P_{U_{t}}$ be a partition
on $U_{t}$ with sdepth $d+2$ and let $[b_{i},c_{i}]$ be the disjoint intervals
starting with $b_{i}=x_{t}f_{i}$, $i\in[e]$. We may suppose that $c_{i}\in C$
for $i\in[e]$. We have $c_{i}=x_{t}w_{ik_{i}}$ for some $1\leq k_{i}\leq r$,
$k_{i}\not=i$ because $C\subset(W)$. Note that $x_{t}f_{k_{i}}\in B$ and so
$k_{i}\leq e$. We consider the intervals $[f_{i},c_{i}]$. These intervals
contain $x_{t}f_{i}$ and $w_{ik_{i}}$. If $w_{ik_{i}}=w_{jk_{j}}$ for
$i\not=j$ then we get $c_{i}=c_{j}$ which is false. Thus these intervals are
disjoint.
Let $I_{e}$ be the ideal generated by $f_{j}$ for $e<j\leq r$ and
$B\setminus(\cup_{i=1}^{e}[f_{i},c_{i}])$. Set $J_{e}=I_{e}\cap J$ and
$U_{e}=I_{e}/J_{e}$. Note that $c_{i}\not\in I_{e}$ for any $i\in[e]$. In the
following exact sequence
$0\rightarrow I_{e}/J_{e}\rightarrow I/J\rightarrow I/J+I_{e}\rightarrow 0$
the last term has a partition of sdepth $d+2$ given by the intervals
$[f_{i},c_{i}]$ for $1\leq i\leq e$. It follows that $I_{e}\not=J_{e}$ because
$\operatorname{sdepth}_{S}I/J=d+1$. Then
$\operatorname{sdepth}_{S}I_{e}/J_{e}\leq d+1$ using [18, Lemma 2.2] and so
$\operatorname{depth}_{S}I_{e}/J_{e}\leq d+1$ by Conjecture 1 applied for
$r^{\prime}<r$. But the last term of the above sequence has depth $>d$ because
$x_{t}$ does not annihilate $f_{i}$ for $i\in[e]$. With the Depth Lemma we get
$\operatorname{depth}_{S}I/J\leq d+1$.
###### Example 4.
Let $n=5$, $r=4$, $f_{1}=x_{2}x_{3}$, $f_{2}=x_{1}x_{2}$, $f_{3}=x_{3}x_{4}$,
$f_{4}=x_{3}x_{5}$ and $J=(x_{1}x_{2}x_{4}x_{5})$. We have
$w_{12}=x_{1}x_{2}x_{3}$, $w_{13}=x_{2}x_{3}x_{4}$, $w_{14}=x_{2}x_{3}x_{5}$,
$w_{34}=x_{3}x_{4}x_{5}$, $w_{23}=x_{1}x_{2}x_{3}x_{4}$,
$w_{24}=x_{1}x_{2}x_{3}x_{5}$, $C_{2}=\\{w_{23},w_{24}\\}$,
$C=C_{2}\cup\\{x_{1}w_{34},x_{2}x_{3}x_{4}x_{5}\\}$,
${\tilde{I}}_{1}=\\{x_{1}f_{3},x_{1}f_{4},f_{2}\\}\supset J$,
${\tilde{I}}_{4}=\\{f_{3},x_{4}f_{2}\\}\supset J$ and
$B\cap(x_{1})=\\{x_{1}f_{3},x_{1}f_{4},x_{4}f_{2},x_{5}f_{2},w_{12}\\}$,
$B\cap(x_{4})=\\{w_{13},w_{14},w_{34},x_{4}f_{2},x_{1}f_{3}\\}$. Note that
$\operatorname{sdepth}_{S}U_{1}\leq d+1=3$,
$\operatorname{sdepth}_{S}U_{4}\leq 3$ because
$|B\cap(x_{1})|=|B\cap(x_{4})|=5>|C\cap(x_{1})|+1=|C\cap(x_{4})|+1=4$. Thus
$\operatorname{depth}_{S}U_{1}=\operatorname{depth}_{S}U_{4}\leq 3$ and so we
get $\operatorname{depth}_{S}I/J\leq 3$ using two different $t$.
###### Theorem 8.
Suppose that $I\subset S$ is minimally generated by four squarefree monomials
$\\{f_{1},\ldots,f_{4}\\}$ of degrees $d$ such that there exists $c\in C$ such
that $\operatorname{supp}c\not\subset\cup_{i\in[4]}\operatorname{supp}f_{i}$.
If $\operatorname{sdepth}_{S}I/J=d+1$ then $\operatorname{depth}_{S}I/J\leq
d+1$.
###### Proof.
Apply Theorem 7, since Conjecture 1 holds for $r<4$ by Theorem 6.
## 3\. Proof of Theorem 6
Suppose that $E\not=\emptyset$ and $s\leq q+r$. For $b=f_{1}x_{i}\in B$ set
$I_{b}=(f_{2},\ldots,f_{r},B\setminus\\{b\\})$, $J_{b}=J\cap I_{b}$. If
$\operatorname{sdepth}_{S}I_{b}/J_{b}\geq d+2$ then let $P_{b}$ be a partition
on $I_{b}/J_{b}$ with sdepth $d+2$. We may choose $P_{b}$ such that each
interval starting with a squarefree monomial of degree $d$, $d+1$ ends with a
monomial of $C$. In $P_{b}$ we have some intervals
$[f_{k},f_{k}x_{i_{k}}x_{j_{k}}]$, $1<k\leq r$ and for all $b^{\prime}\in
B\setminus\\{b,f_{2}x_{i_{2}},f_{2}x_{j_{2}},\ldots,f_{r}x_{i_{r}},f_{r}x_{j_{r}}\\}]$
an interval $[b^{\prime},c_{b^{\prime}}]$. We define
$h:[[\\{f_{2},\ldots,f_{r}\\}\cup
B]\setminus\\{b,f_{2}x_{i_{2}},f_{2}x_{j_{2}},\ldots,f_{r}x_{i_{r}},f_{r}x_{j_{r}}\\}]\rightarrow
C$ by $f_{k}\rightarrow f_{k}x_{i_{k}}x_{j_{k}}$ and $b^{\prime}\rightarrow
c_{b^{\prime}}$. Then $h$ is an injection and $|\operatorname{Im}\ h|=s-r\leq
q$ (if $s=r+q$ then $h$ is a bijection).
###### Lemma 5.
Suppose that the following conditions hold:
1. (1)
$r=2$, $4\leq s\leq q+2$,
2. (2)
$C\subset((f_{1})\cap(f_{2}))\cup((E)\cap(f_{1},f_{2}))\cup(\cup_{a,a^{\prime}\in
E,a\not=a^{\prime}}(a)\cap(a^{\prime}))$,
3. (3)
$\operatorname{sdepth}_{S}I_{b}/J_{b}\geq d+2$ for a
$b\in(B\cap(f_{1}))\setminus(f_{2})$.
Then either $\operatorname{sdepth}_{S}I/J\geq d+2$, or there exists a nonzero
ideal $I^{\prime}\subsetneq I$ generated by a subset of $\\{f_{1},f_{2}\\}\cup
B$ such that $\operatorname{sdepth}_{S}I^{\prime}/J^{\prime}\leq d+1$ for
$J^{\prime}=J\cap I^{\prime}$ and
$\operatorname{depth}_{S}I/(J,I^{\prime})\geq d+1$.
###### Proof.
Since $\operatorname{sdepth}_{S}I_{b}/J_{b}\geq d+2$ we consider $h$ as above
for a partition $P_{b}$ with sdepth $d+2$ of $I_{b}/J_{b}$. We have an
interval $[f_{2},c^{\prime}_{2}]$ in $P_{b}$. Suppose that
$B\cap[f_{2},c^{\prime}_{2}]=\\{u,u^{\prime}\\}$. A sequence
$a_{1},\ldots,a_{k}$ is called a path from $a_{1}$ to $a_{k}$ if the following
statements hold:
(i) $a_{l}\in B\setminus\\{b,u,u^{\prime}\\}$, $l\in[k]$,
(ii) $a_{l}\not=a_{j}$ for $1\leq l<j\leq k$,
(iii) $a_{l+1}|h(a_{l})$ and $h(a_{l})\not\in(b)$ for all $1\leq l<k$.
This path is weak if $h(a_{j})\in(u,u^{\prime})$ for some $j\in[k]$. It is bad
if $h(a_{k})\in(b)$ and it is maximal if either $h(a_{k})\in(b)$, or all
divisors from $B$ of $h(a_{k})$ are in
$\\{u,u^{\prime},a_{1},\ldots,a_{k}\\}$. If $a=a_{1}$ we say that the above
path starts with $a$.
By hypothesis $s\geq 4$ and so there exists $a_{1}\in
B\setminus\\{b,u,u^{\prime}\\}$. Set $c_{1}=h(a_{1})$. If $c_{1}\in(b)$ then
the path $\\{a_{1}\\}$ is maximal and bad. We construct below, as an example,
a path with $k>1$. By recurrence choose if possible $a_{p+1}$ to be a divisor
from $B$ of $c_{p}$, which is not in $\\{b,u,u^{\prime},a_{1},\ldots,a_{p}\\}$
and set $c_{p}=h(a_{p})$, $p\geq 1$. This construction ends at step $p=e$ if
all divisors from $B$ of $c_{e}$ are in
$\\{b,u,u^{\prime},a_{1},\ldots,a_{e}\\}$. If $c_{i}\not\in(b)$ for $1\leq
i<e$ then $\\{a_{1},\ldots,a_{e}\\}$ is a maximal path. If one
$c_{p}\in(u,u^{\prime})$ then the constructed path is weak. If $c_{e}\in(b)$
then this path is bad. We have three cases:
1) there exist no weak path and no bad path starting with $a_{1}$,
2) there exists a weak path starting with $a_{1}$ but no bad path starts with
$a_{1}$,
3) there exists a bad path starting with $a_{1}$.
In the first case, set $T_{1}=\\{b^{\prime}\in B:\mbox{there\ \ exists\ \ a\ \
path}\ \ a_{1},\ldots,a_{k}\ \ \mbox{with}\ \ a_{k}=b^{\prime}\\}$,
$G_{1}=B\setminus T_{1}$, and $I^{\prime}_{1}=(f_{1},G_{1})$,
$I^{\prime}_{2}=(f_{2},G_{1})$, $I^{\prime}_{12}=(f_{1},f_{2},G_{1})$,
$I^{\prime\prime}=(G_{1})$, $J^{\prime}_{1}=I^{\prime}_{1}\cap J$,
$J^{\prime}_{2}=I^{\prime}_{2}\cap J$, $J^{\prime}_{12}=I^{\prime}_{12}\cap
J$, $J^{\prime\prime}=I^{\prime\prime}\cap J$. Note that
$I^{\prime\prime}\not=0$ because $b\in I^{\prime\prime}$ and all divisors from
$B$ of a monomial $c\in U_{1}=h(T_{1})$ belong to $T_{1}$. Consider the
following exact sequence
$0\rightarrow I^{\prime}_{12}/J^{\prime}_{12}\rightarrow I/J\rightarrow
I/(J,I^{\prime}_{12})\rightarrow 0.$
If $U_{1}\cap(f_{1},f_{2})=\emptyset$ then the last term has depth $\geq d+1$
and sdepth $\geq d+2$ using the restriction of $P_{b}$ to
$(T_{1})\setminus(J,I^{\prime}_{12})$ since $h(b^{\prime})\notin
I^{\prime}_{12}$, for all $b^{\prime}\in T_{1}$. When the first term has
sdepth $\geq d+2$ then by [18, Lemma 2.2] the middle term has sdepth $\geq
d+2$. Otherwise, the first term has sdepth $\leq d+1$ and we may take
$I^{\prime}=I^{\prime}_{12}$.
If $U_{1}\cap(f_{1})=\emptyset$, but $b_{2}\in T_{1}\cap(f_{2})$, then in the
following exact sequence
$0\rightarrow I^{\prime}_{1}/J^{\prime}_{1}\rightarrow I/J\rightarrow
I/(J,I^{\prime}_{1})\rightarrow 0$
the last term has sdepth $\geq d+2$ since $h(b^{\prime})\notin
I^{\prime}_{1}$, for all $b^{\prime}\in T_{1}$ and we may substitute the
interval $[b_{2},h(b_{2})]$ from the restriction of $P_{b}$ by
$[f_{2},h(b_{2})]$, the second monomial from $[f_{2},h(b_{2})]\cap B$ being
also in $T_{1}$. As above we get either $\operatorname{sdepth}_{S}I/J\geq
d+2$, or $\operatorname{sdepth}_{S}I^{\prime}_{1}/J^{\prime}_{1}\leq d+1$,
$\operatorname{depth}_{S}I/(J,I^{\prime}_{1})\geq d+1$. Similarly, we do when
$U_{1}\cap(f_{2})=\emptyset$ but $U_{1}\cap(f_{1})\not=\emptyset$.
Now, suppose that $b_{1}\in T_{1}\cap(f_{1})$ and $b_{2}\in T_{1}\cap(f_{2})$.
We claim to choose $b_{1}\not=b_{2}$ and such that one from
$h(b_{1}),h(b_{2})$ is not in $(w_{12})$, let us say
$h(b_{1})\not\in(w_{12})$. Indeed, if $w_{12}\not\in B$ and
$h(b_{1}),h(b_{2})\in(w_{12})$ then necessarily $h(b_{1})=h(b_{2})$ and it
follows $b_{1}=b_{2}=w_{12}$ which is false. Suppose that $w_{12}\in B$ and
$h(b_{2})=x_{j}w_{12}$. Then choose $b_{1}=x_{j}f_{1}\in T_{1}$. If
$h(b_{1})\in(w_{12})$ then we get $h(b_{1})=h(b_{2})$ and so
$b_{1}=b_{2}=w_{12}$ which is impossible.
In the following exact sequence
$0\rightarrow I^{\prime\prime}/J^{\prime\prime}\rightarrow I/J\rightarrow
I/(J,I^{\prime\prime})\rightarrow 0$
the last term has sdepth $\geq d+2$ since we may replace the intervals
$[b_{1},h(b_{1})]$, $[b_{2},h(b_{2})]$ of the restriction of $P_{b}$ to
$(T_{1})\setminus(J,I^{\prime\prime})$ with the disjoint intervals
$[f_{1},h(b_{1})]$, $[f_{2},h(b_{2})]$. Also the last term has depth $\geq
d+1$ because in the exact sequence
$0\rightarrow(f_{2})/(J,I^{\prime\prime})\cap(f_{2})\rightarrow
I/(J,I^{\prime\prime})\rightarrow I/(J,I^{\prime\prime},f_{2})\rightarrow 0$
the end terms have depth $\geq d+1$ since $h(b_{1})\not\in(f_{2})$, otherwise
$h(b_{1})\in(w_{12})$, which is false. As above we get either
$\operatorname{sdepth}_{S}I/J\geq d+2$, or
$\operatorname{sdepth}_{S}I^{\prime\prime}/J^{\prime\prime}\leq d+1$,
$\operatorname{depth}_{S}I/(J,I^{\prime\prime})\geq d+1$.
In the second case, let $a_{1},\ldots,a_{t_{1}}$ be a weak path and set
$c_{j}=h(a_{j})$ for $j\in[t_{1}]$. We may suppose that $c_{t_{1}}\in(u)$,
otherwise take a shorter path. Denote $T_{1}$, $U_{1}$ as in the first case,
which we keep it fix even we will change a little $h$. Suppose that
$a_{t_{1}}\in(f_{2})$. Then change in $P_{b}$ the intervals
$[a_{t_{1}},c_{t_{1}}]$, $[f_{2},c^{\prime}_{2}]$ by $[f_{2},c_{t_{1}}]$,
$[u^{\prime},c^{\prime}_{2}]$. Thus the new $c^{\prime}_{2}$ is among
$\\{c_{1},\ldots,c_{t_{1}}\\}\subset U_{1}$, though the old
$c^{\prime}_{2}\not\in U_{1}$. Also the new $u^{\prime}$ is in $T_{1}$.
However, if the old $u^{\prime}$ is not a divisor of a $c$ from $U_{1}$, then
the proof goes as in the first case using $T^{\prime}_{1}=T_{1}\cup\\{u\\}$,
$G^{\prime}_{1}=B\setminus T^{\prime}_{1}$ with $I^{\prime}=I^{\prime}_{2}$,
or $I^{\prime}=I^{\prime}_{12}$. Otherwise, $T^{\prime}_{1}$ should be
completed because $u^{\prime}$ is not now in $[f_{2},c^{\prime}_{2}]$ and we
may consider some paths starting with $u^{\prime}$. Note that there exists a
path from $a_{1}$ to $u^{\prime}$ since $u^{\prime}$ is a divisor of a
monomial from $U_{1}$. It follows that there exist no bad path starting with
$u^{\prime}$. Take ${\tilde{T}_{1}}=T^{\prime}_{1}\cup\\{b^{\prime}\in
B:\mbox{there\ \ exists\ \ a\ \ path\ \ from}\ \ u^{\prime}\ \ \mbox{to}\ \
b^{\prime}\\}$ and the proof goes as above with $\tilde{T}_{1}$ instead
$T^{\prime}_{1}$, that is with $I^{\prime}$ generated by a subset of
$\\{f_{1},f_{2}\\}\cup{\tilde{G}}_{1}$ for
${\tilde{G}}_{1}=B\setminus{\tilde{T}}_{1}$.
Now suppose that $a_{t_{1}}\not\in(f_{2})$ but there exists $1\leq v<t_{1}$
such that $a_{v}\in(f_{2})$ and $a_{v}|c_{t_{1}}$. Then we may replace in
$P_{b}$ the intervals $[a_{p},c_{p}],v\leq p\leq t_{1}$ with the intervals
$[a_{v},c_{t_{1}}],[a_{p+1},c_{p}],v\leq p<t_{1}$. The old $c_{t_{1}}$ becomes
the new $c_{v}$, that is we reduce to the case when $u$ divides $c_{v}$ and
$a_{v}\in(f_{2})$, subcase solved above.
Remains to study the subcase when there exist no $a_{v}\in(f_{2})$, $1\leq
v\leq t_{1}$ with $a_{v}|c_{t_{1}}$. Then there exists an $a_{t_{1}+1}\in
B\cap(f_{2})$, $a_{t_{1}+1}\not=u$ such that $a_{t_{1}+1}|c_{t_{1}}$. Clearly,
$a_{t_{1}+1}\not=u^{\prime}$ because otherwise $c^{\prime}_{2}=c_{t_{1}}$. We
have two subcases:
$1^{\prime})$ there exists a path $a_{t_{1}+1},\ldots,a_{l}$ such that
$h(a_{l})\in(a_{v^{\prime}})$ for some $1\leq v^{\prime}\leq t_{1}$,
$2^{\prime})$ for any path $a_{t_{1}+1},\ldots,a_{p}$, any $h(a_{j})$,
$t_{1}<j\leq p$ does not belong to $(a_{1},\ldots,a_{t_{1}})$.
In the first subcase, we replace in $P_{b}$ the intervals
$[a_{j},c_{j}],v^{\prime}\leq j\leq l$ with the intervals
$[a_{v^{\prime}},c_{l}],[a_{j+1},c_{j}],v^{\prime}\leq j<l$. The new
$h(a_{t_{1}+1})$ is the old $c_{t_{1}}$ and we may proceed as above. In the
second case, we set
$T_{2}=\\{b^{\prime}\in B:\mbox{there\ \ exists\ \ a\ \ path\ \ from}\ \
a_{t_{1}+1}\ \ \mbox{to}\ \ b^{\prime}\\}.$
Note that any path starting from $a_{t_{1}+1}$ can be completed to a path from
$a_{1}$ by adding the monomials $a_{1},\ldots,a_{t_{1}}$. Thus there exists no
bad path starting with $a_{t_{1}+1}$, otherwise we can get one starting from
$a_{1}$, which is false.
If there exists no weak path starting with $a_{t_{1}+1}$ then we proceed as in
the first case with $T_{2}$ instead $T_{1}$. If there exists a weak path
starting with $a_{t_{1}+1}$ then we proceed as above in case 2) with
$T^{\prime}_{2}$, or ${\tilde{T}}_{2}$ instead $T^{\prime}_{1}$, or
${\tilde{T}}_{1}$, except in the subcase $2^{\prime})$ when we will define
similarly a $T_{3}$ given by the paths starting with a certain $a_{t_{2}+1}$.
Note that the whole set $\\{a_{1},\ldots,a_{t_{2}}\\}$ has different
monomials. After several such steps we must arrive in the case $p=t_{m}$ when
$\\{a_{1},\ldots,a_{t_{m}}\\}$ has different monomials and the subcase
$2^{\prime})$ does not appear. We end this case using $T_{m}$, or
$T^{\prime}_{m}$, or ${\tilde{T}}_{m}$ instead $T_{1}$, or $T^{\prime}_{1}$,
or ${\tilde{T}}_{1}$.
In the third case, let $a_{1},\ldots,a_{t_{1}}$ be a bad path starting with
$a_{1}$. Set $c_{j}=h(a_{j})$, $j\in[t_{1}]$. Then $c_{t_{1}}=bx_{l_{1}}$ and
let us say $b=f_{1}x_{i}$. If $a_{t_{1}}\in(f_{1})$ then changing in $P_{b}$
the interval $[a_{t_{1}},c_{t_{1}}]$ by $[f_{1},c_{t_{1}}]$ we get a partition
on $I/J$ with sdepth $d+2$. Thus we may assume that $a_{t_{1}}\not\in(f_{1})$.
If $f_{1}x_{l_{1}}\in\\{a_{1},\ldots,a_{t_{1}-1}\\}$, let us say
$fx_{l_{1}}=a_{v}$, $1\leq v<t_{1}$ then we may replace in $P_{b}$ the
intervals $[a_{p},c_{p}],v\leq p\leq t_{1}$ with the intervals
$[a_{v},c_{t_{1}}],[a_{p+1},c_{p}],v\leq p<t_{1}$. Now we see that we have in
$P_{b}$ the interval $[f_{1}x_{l_{1}},f_{1}x_{i}x_{l_{1}}]$ and switching it
with the interval $[f_{1},f_{1}x_{i}x_{l_{1}}]$ we get a partition with sdepth
$\geq d+2$ for $I/J$.
Thus we may assume that $f_{1}x_{l_{1}}\notin\\{a_{1},...,a_{t_{1}}\\}$. Now
set $a_{t_{1}+1}=fx_{l_{1}}$. Let
$a_{t_{1}+1},\ldots,a_{k}$ be a path starting with $a_{t_{1}+1}$ and set
$c_{j}=h(a_{j})$, $t_{1}<j\leq k$. If $a_{p}=a_{v}$ for $v\leq t_{1}$,
$p>t_{1}$ then change in $P_{b}$ the intervals $[a_{j},c_{j}],v\leq j\leq p$
with the intervals $[a_{v},c_{p}],[a_{j+1},c_{j}],v\leq j<p$. We have in
$P_{b}$ an interval $[f_{1}x_{l_{1}},f_{1}x_{i}x_{l_{1}}]$ and switching it to
$[f_{1},f_{1}x_{i}x_{l_{1}}]$ we get a partition with sdepth $\geq d+2$ for
$I/J$. Thus we may suppose that in fact
$a_{p}\not\in\\{b,a_{1},\ldots,a_{p-1}\\}$ for any $p>t_{1}$ (with respect to
any path starting with $a_{t_{1}+1}$). We have three subcases:
$1^{\prime\prime})$ there exist no weak path and no bad path starting with
$a_{t_{1}+1}$,
$2^{\prime\prime})$ there exists a weak path starting with $a_{t_{1}+1}$ but
no bad path starts with $a_{t_{1}+1}$,
$3^{\prime\prime})$ there exists a bad path starting with $a_{t_{1}+1}$.
Set $T_{2}=\\{b^{\prime}\in B:\mbox{there\ exists\ a\ path}\ \
a_{t_{1}+1},\ldots,a_{k}\ \ \mbox{with}\ \ a_{k}=b^{\prime}\\}$,
$G_{2}=B\setminus T_{2}$, $U_{2}=h(T_{2})$ in the first subcase, and see that
$I^{\prime}$ generated by a subset of $\\{f_{1},f_{2}\\}\cup G_{2}$ chosen as
above works.
In the second subcase, let $a_{t_{1}+1},\ldots,a_{k}$ be a weak path and set
$c_{j}=h(a_{j})$ for $t_{1}<j\leq k$. We may suppose that $c_{k}\in(u)$.
Changing $P_{b}$ we may suppose that the new $c^{\prime}_{2}$ is in $U_{2}$ as
above. If the old $u^{\prime}$ was not a divisor of a $c\in U_{2}$ then the
proof goes as in the first case with $T^{\prime}_{2}=T_{2}\cup\\{u\\}$,
$I^{\prime}=I^{\prime}_{2}$. Otherwise, $T^{\prime}_{2}$ should be completed
to a ${\tilde{T}_{2}}$ similar to ${\tilde{T}_{1}}$. The proof goes as above
with $\tilde{T}_{2}$ instead $T^{\prime}_{2}$.
In the third subcase, let $a_{t_{1}+1},\ldots,a_{t_{2}}$ be a bad path
starting with $a_{t_{1}+1}$ and set $c_{j}=h(a_{j})$ for $j>t_{1}$. We saw
that the whole set $\\{a_{1},\ldots,a_{t_{2}}\\}$ has different monomials. As
above $c_{t_{2}}=bx_{l_{2}}$ and we may reduce to the case when
$f_{1}x_{l_{2}}\not\in\\{a_{1},\ldots,a_{t_{1}}\\}$. Set
$a_{t_{2}+1}=f_{1}x_{l_{2}}$ and again we consider three subcases, which we
treat as above. Anyway after several such steps we must arrive in the case
$p=t_{m}$ when $b|c_{t_{m}}$ and again a certain $f_{1}x_{l_{m}}$ is not among
$\\{a_{1},\ldots,a_{t_{m}}\\}$ and there exist no bad path starting with
$a_{t_{m}+1}=f_{1}x_{l_{m}}$. This follows since we may reduce to the case
when the set $\\{a_{1},\ldots,a_{t_{m}}\\}$ has different monomials and so the
procedures should stop for some $m$. Finally, using
$T_{m}=\\{b^{\prime}\in B:\mbox{there\ exists\ a\ path}\ \
a_{t_{m}+1},\ldots,a_{k}\ \ \mbox{with}\ \ a_{k}=b^{\prime}\\}$
(resp. $T^{\prime}_{m}$, or ${\tilde{T}}_{m}$) as $T_{1}$ (resp.
$T^{\prime}_{1}$, or ${\tilde{T}}_{1}$) above we are done.
###### Lemma 6.
Suppose that the following conditions hold:
1. (1)
$r=3$, $6\leq s\leq q+3$,
2. (2)
$C\subset(\cup_{i,j\in[3],i\not=j}(f_{i})\cap(f_{j}))\cup((E)\cap(f_{1},f_{2},f_{3}))\cup(\cup_{a,a^{\prime}\in
E,a\not=a^{\prime}}(a)\cap(a^{\prime}))$,
3. (3)
There exists $b\in(B\cap(f_{1}))\setminus(f_{2},f_{3})$ such that
$\operatorname{sdepth}_{S}I_{b}/J_{b}\geq d+2$.
Then either $\operatorname{sdepth}_{S}I/J\geq d+2$, or there exists a nonzero
ideal $I^{\prime}\subsetneq I$ generated by a subset of
$\\{f_{1},f_{2},f_{3}\\}\cup B$ such that
$\operatorname{sdepth}_{S}I^{\prime}/J^{\prime}\leq d+1$ for $J^{\prime}=J\cap
I^{\prime}$ and $\operatorname{depth}_{S}I/(J,I^{\prime})\geq d+1$.
###### Proof.
We follow the proof of Lemma 5. Since
$\operatorname{sdepth}_{S}I_{b}/J_{b}\geq d+2$ we consider $h$ as above for a
partition $P_{b}$ with sdepth $d+2$ of $I_{b}/J_{b}$. We have two intervals
$[f_{2},c^{\prime}_{2}]$, $[f_{3},c^{\prime}_{3}]$ in $P_{b}$. Suppose that
$B\cap[f_{i},c^{\prime}_{i}]=\\{u_{i},u^{\prime}_{i}\\}$, $1<i\leq 3$ . As in
Lemma 5 we define a path $a_{1},\ldots,a_{k}$ from $a_{1}$ to $a_{k}$ and a
bad path. The above path is weak if
$h(a_{j})\in(u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3})$ for some $j\in[k]$.
It is maximal if either $h(a_{k})\in(b)$, or all divisors from $B$ of
$h(a_{k})$ are in
$\\{b,u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3},a_{1},\ldots,a_{k}\\}$.
By hypothesis $s\geq 6$ and there exists $a_{1}\in
B\setminus\\{b,u_{2},u_{2}^{\prime},u_{3},u^{\prime}_{3}\\}$. Set
$c_{1}=h(a_{1})$. If $c_{1}\in(b)$ then the path $\\{a_{1}\\}$ is maximal and
bad. We construct below a path with $k>1$. By recurrence choose if possible
$a_{p+1}$ to be a divisor from $B$ of $c_{p}$, which is not in
$\\{b,u_{2},u^{\prime}_{2},u_{3},u_{3}^{\prime},a_{1},\ldots,a_{p}\\}$ and set
$c_{p}=h(a_{p})$, $p\geq 1$. This construction ends at step $p=e$ if all
divisors from $B$ of $c_{e}$ are in
$\\{b,u_{2},u^{\prime}_{2},u_{3},u_{3}^{\prime},a_{1},\ldots,a_{e}\\}$. If
$c_{j}\not\in(b)$ for $1\leq j<e$ then $\\{a_{1},\ldots,a_{e}\\}$ is a maximal
path. If one $c_{p}\in(u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3})$ then the
constructed path is weak. If $c_{e}\in(b)$ then this path is bad.
We may reduce to the situation when $P_{b}$ satisfies the following property:
$(*)$ For all $1\leq i<j\leq 3$ if $w_{ij}\in
B\setminus\\{u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3}\\}$, $1\leq i<j\leq 3$
then $h(w_{ij})\not\in(u_{i},u^{\prime}_{i},u_{j},u^{\prime}_{j})$ if $i>1$.
Indeed, suppose that $w_{ij}\in
B\setminus\\{u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3}\\}$ and
$h(w_{ij})\in(u_{j})$. Then $h(w_{ij})=x_{l}w_{ij}$ for some
$l\not\in\operatorname{supp}w_{ij}$ and we must have let us say
$u_{j}=x_{l}f_{j}$. Changing in $P_{b}$ the intervals
$[f_{j},c^{\prime}_{j}]$, $[w_{ij},h(w_{ij})]$ with $[f_{j},h(w_{ij})]$,
$[u^{\prime}_{j},c^{\prime}_{j}]$ we see that we may assume $u_{j}=w_{ij}$.
Suppose that $(*)$ holds. We have three cases:
1) there exist no weak path and no bad path starting with $a_{1}$,
2) there exists a weak path starting with $a_{1}$ but no bad path starts with
$a_{1}$,
3) there exists a bad path starting with $a_{1}$.
In the first case, set $T_{1}=\\{b^{\prime}\in B:\mbox{there\ \ exists\ \ a\ \
path}\ \ a_{1},\ldots,a_{t}\ \ \mbox{with}\ \ a_{t}=b^{\prime}\\}$,
$G_{1}=B\setminus T_{1}$, and for $k=(k_{1},\ldots,k_{m})$, $1\leq
k_{1}<\ldots k_{m}\leq 3$, $0\leq m\leq 3$ set
$I^{\prime}_{k}=(f_{k_{1}},\ldots,f_{k_{m}},G_{1})$,
$J^{\prime}_{k}=I^{\prime}_{k}\cap J$, and $I^{\prime}_{0}=(G_{1})$,
$J^{\prime}_{0}=I^{\prime}_{0}\cap J$ for $m=0$. Note that all divisors from
$B$ of a monomial $c\in U_{1}=h(T_{1})$ belong to $T_{1}$, and
$I^{\prime}_{0}\not=0$ because $b\in I^{\prime}_{0}$. Consider the following
exact sequence
$0\rightarrow I^{\prime}_{k}/J^{\prime}_{k}\rightarrow I/J\rightarrow
I/(J,I^{\prime}_{k})\rightarrow 0.$
If $U_{1}\cap(f_{1},f_{2},f_{3})=\emptyset$ then the last term of the above
exact sequence given for $k=(1,2,3)$ has depth $\geq d+1$ and sdepth $\geq
d+2$ using the restriction of $P_{b}$ to $(T_{1})\setminus(J,I^{\prime}_{k})$
since $h(b^{\prime})\notin I^{\prime}_{k}$ , for all $b^{\prime}\in T_{1}$.
When the first term has sdepth $\geq d+2$ then by [18, Lemma 2.2] the middle
term has sdepth $\geq d+2$ which is enough.
If $U_{1}\cap(f_{1},f_{2})=\emptyset$, but there exists $b_{3}\in
T_{1}\cap(f_{3})$, then set $k=(1,2)$. In the following exact sequence
$0\rightarrow I^{\prime}_{k}/J^{\prime}_{k}\rightarrow I/J\rightarrow
I/(J,I^{\prime}_{k})\rightarrow 0$
the last term has sdepth $\geq d+2$ since $h(b^{\prime})\notin
I^{\prime}_{k}$, for all $b^{\prime}\in T_{1}$ and we may substitute the
interval $[b_{3},h(b_{3})]$ from the restriction of $P_{b}$ by
$[f_{3},h(b_{3})]$, the second monomial from $[f_{3},h(b_{3})]\cap B$ being
also in $T_{1}$. As above we get either $\operatorname{sdepth}_{S}I/J\geq
d+2$, or $\operatorname{sdepth}_{S}I^{\prime}_{1}/J^{\prime}_{1}\leq d+1$,
$\operatorname{depth}_{S}I/(J,I^{\prime}_{1})\geq d+1$.
Now, we omit other subcases considering only the worst subcase $m=0$. Let
$b_{1}\in T_{1}\cap(f_{1})$, $b_{2}\in T_{1}\cap(f_{2})$ and $b_{3}\in
T_{1}\cap(f_{3})$. For $1\leq l<j\leq 3$ we claim that we may choose
$b_{l}\not=b_{j}$ and such that one from $h(b_{l}),h(b_{j})$ is not in
$(w_{lj})$. Indeed, if $w_{lj}\not\in B$ and $h(b_{l}),h(b_{j})\in(w_{lj})$
then necessarily $h(b_{l})=h(b_{j})$ and it follows $b_{l}=b_{j}=w_{lj}$,
which is false. Suppose that $w_{lj}\in B$ and $h(b_{j})=x_{p}w_{lj}$. Then
choose $b_{l}=x_{p}f_{l}\in T_{1}$. If $h(b_{l})\in(w_{lj})$ then we get
$h(b_{l})=h(b_{j})$ and so $b_{l}=b_{j}=w_{lj}$, which is impossible.
Therefore, we may choose $b_{j}$ such that $h(b_{1})\not\in(w_{12})$,
$h(b_{2})\not\in(w_{23})$. Note that it is possible that $f_{1}|c$ for some
$c\in h(T_{1})$ even $b\not|c$ for any $c\in U_{1}$. If $h(b_{1})\in(w_{13})$
then we may also choose $h(b_{3})\not\in(w_{13})$. In the case when
$h(b_{1})\not\in(w_{13})$, choose any $b_{3}\in T_{1}\cap(f_{3})$ different
from the others $b_{j}$. We conclude that the possible intervals
$[f_{j},h(b_{j})]$, $j\in[3]$ are disjoint. Next we change the intervals
$[b_{j},h(b_{j})]$, $j\in[3]$ from the restriction of $P_{b}$ to
$(T_{1})\setminus(J,I^{\prime}_{0})$ by $[f_{j},h(b_{j})]$, the second
monomial from $[f_{j},h(b_{j})]\cap B$ being also in $T_{1}$. We claim that
$I/(J,I^{\prime}_{0})$ has depth $\geq d+1$. Indeed, in the following exact
sequence
$0\rightarrow(f_{2})/(f_{2})\cap(J,I^{\prime}_{0},f_{3})\rightarrow
I/(J,I^{\prime}_{0},f_{3})\rightarrow
I/(J,I^{\prime}_{0},f_{2},f_{3})\rightarrow 0$
the first term has depth $\geq d+1$ because
$h(b_{2})\not\in(f_{2})\cap(f_{3})$. If $h(b_{1})\not\in(f_{3})$ then
$h(b_{1})\not\in(f_{2},f_{3})\cap(f_{1})$ and so the last term has depth $\geq
d+1$. If $h(b_{1})\in(w_{13})$ then we may find a
$b^{\prime}\in(B\cap(f_{1}))\setminus(f_{3})$ dividing $h(b_{1})$. It follows
that $b^{\prime}\in T_{1}$ and $b^{\prime}\not\in(f_{2},f_{3})\cap(f_{1})$,
which implies that the last term has again depth $\geq d+1$. Thus
$\operatorname{depth}_{S}I/(J,I^{\prime}_{0},f_{3})\geq d+1$ by the Depth
Lemma. Our claim follows from the exact sequence
$0\rightarrow(f_{3})/(f_{3})\cap(J,I^{\prime}_{0})\rightarrow
I/(J,I^{\prime}_{0})\rightarrow I/(J,I^{\prime}_{0},f_{3})\rightarrow 0$
because the first term has depth $\geq d+1$. Therefore, as above we get either
$\operatorname{sdepth}_{S}I/J\geq d+2$, or
$\operatorname{sdepth}_{S}I^{\prime}_{0}/J^{\prime}_{0}\leq d+1$,
$\operatorname{depth}_{S}I/(J,I^{\prime}_{0})\geq d+1$.
In the second case, let $a_{1},\ldots,a_{t_{1}}$ be a weak path and set
$c_{j}=h(a_{j})$ for $j\in[t_{1}]$. We may suppose that $c_{t_{1}}\in(u_{2})$,
otherwise take a shorter path. Denote $T_{1}$, $U_{1}$ as in the first case.
First consider the subcase when $U_{1}\cap(f_{3})=\emptyset$. Suppose that
$a_{t_{1}}\in(f_{2})$. Then change in $P_{b}$ the intervals
$[a_{t_{1}},c_{t_{1}}]$, $[f_{2},c^{\prime}_{2}]$ by $[f_{2},c_{t_{1}}]$,
$[u^{\prime}_{2},c^{\prime}_{2}]$. Thus the new $c^{\prime}_{2}$ is among
$\\{c_{1},\ldots,c_{t_{1}}\\}\subset U_{1}$, though the old
$c^{\prime}_{2}\not\in U_{1}$. If the old $u^{\prime}_{2}$ is not a divisor of
any $c\in U_{1}$ then the proof goes as in the first case with
$T^{\prime}_{1}=T_{1}\cup\\{u_{2}\\}$. If the old $u^{\prime}_{2}$ is a
divisor of a monomial from $U_{1}$ then $T^{\prime}_{1}$ should be completed
because the old $u^{\prime}_{2}$ is not now in $[f_{2},c^{\prime}_{2}]$. Note
that there exists a path from $a_{1}$ to $u^{\prime}_{2}$ since
$u^{\prime}_{2}$ is a divisor of a monomial from $U_{1}$. It follows that
there exist no bad path starting with $u^{\prime}_{2}$. It is worth to mention
that the old $c^{\prime}_{2}$ is now in $U_{1}$ and we should consider all
pathes starting with divisors of $c^{\prime}_{2}$ from $B$. Take
${\tilde{T}_{1}}=T^{\prime}_{1}\cup\\{b^{\prime}\in B:\mbox{there\ \ exists\ \
a\ \ path\ \ from}\ \ u^{\prime}_{2}\ \ \mbox{to}\ \ b^{\prime}\\}$ and the
proof goes as above with $\tilde{T}_{1}$ instead $T^{\prime}_{1}$, that is
with $I^{\prime}$ generated by a subset of
$\\{f_{1},f_{2},f_{3}\\}\cup{\tilde{G}}_{1}$, where
${\tilde{G}}_{1}=B\setminus{\tilde{T}}_{1}$.
Now suppose that $a_{t_{1}}\not\in(f_{2})$ but there exists $1\leq v<t_{1}$
such that $a_{v}\in(f_{2})$ and $a_{v}|c_{t_{1}}$. Then we may replace in
$P_{b}$ the intervals $[a_{p},c_{p}],v\leq p\leq t_{1}$ with the intervals
$[a_{v},c_{t_{1}}],[a_{p+1},c_{p}],v\leq p<t_{1}$. The old $c_{t_{1}}$ becomes
the new $c_{v}$, that is we reduce to the case when $u_{2}$ divides $c_{v}$
and $a_{v}\in(f_{2})$, subcase solved above.
Remains to study the subcase when there exist no $a_{v}\in(f_{2})$, $1\leq
v\leq t_{1}$ with $a_{v}|c_{t_{1}}$. Then there exists an $a_{t_{1}+1}\in
B\cap(f_{2})$, $a_{t_{1}+1}\not=u_{2}$ such that $a_{t_{1}+1}|c_{t_{1}}$.
Clearly, $a_{t_{1}+1}\not=u^{\prime}_{2}$ because otherwise
$c^{\prime}_{2}=c_{t_{1}}$. We have two subcases:
$1^{\prime})$ there exists a path $a_{t_{1}+1},\ldots,a_{l}$ such that
$h(a_{l})\in(a_{v^{\prime}})$ for some $1\leq v^{\prime}\leq t_{1}$,
$2^{\prime})$ for any path $a_{t_{1}+1},\ldots,a_{p}$, any $h(a_{j})$,
$t_{1}<j\leq p$ does not belong to $(a_{1},\ldots,a_{t_{1}})$.
In the first subcase, we replace in $P_{b}$ the intervals
$[a_{j},c_{j}],v^{\prime}\leq j\leq l$ with the intervals
$[a_{v^{\prime}},c_{l}],[a_{j+1},c_{j}],v^{\prime}\leq j<l$. The new
$h(a_{t_{1}+1})$ is the old $c_{t_{1}}$ and we may proceed as above. In the
second subcase we set
$T_{2}=\\{b^{\prime}\in B:\mbox{there\ \ exists\ \ a\ \ path\ \ from}\ \
a_{t_{1}+1}\ \ \mbox{to}\ \ b^{\prime}\\}.$
Note that any path starting from $a_{t_{1}+1}$ can be completed to a path from
$a_{1}$ by adding the monomials $a_{1},\ldots,a_{t_{1}}$. Thus there exists no
bad path starting with $a_{t_{1}+1}$, otherwise we can get one starting from
$a_{1}$, which is false.
If there exists no weak path starting with $a_{t_{1}+1}$ then we proceed as in
the first case with $T_{2}$ instead $T_{1}$. If there exists a weak path
starting with $a_{t_{1}+1}$ then we proceed as above in case 2) with
$T^{\prime}_{2}$, or ${\tilde{T}}_{2}$ instead $T^{\prime}_{1}$, or
${\tilde{T}}_{1}$, except in the subcase $2^{\prime})$ when we will define
similarly a $T_{3}$ given by the paths starting with a certain $a_{t_{2}+1}$.
Note that the whole set $\\{a_{1},\ldots,a_{t_{2}}\\}$ has different
monomials. After several such steps we must arrive in the case $p=t_{m}$ when
$\\{a_{1},\ldots,a_{t_{m}}\\}$ has different monomials and the subcase
$2^{\prime})$ does not appear. We end this case using $T_{m}$, or
$T^{\prime}_{m}$, or ${\tilde{T}}_{m}$ instead $T_{1}$, $T^{\prime}_{1}$, or
${\tilde{T}}_{1}$. We should mention that if there exists $b_{1}\in T_{m}$ (or
in $T^{\prime}_{m}$, ${\tilde{T}}_{m}$) such that $h(b_{1})\in(f_{1})$ then
changing $P_{b}$ as in case 1) we may suppose that $h(b_{1})\not\in(w_{12})$
and $b_{1}\in T_{m}\cap(f_{1})$. Thus we may consider the interval
$[f_{1},h(b_{1})]$ disjoint of $[f_{2},c^{\prime}_{2}]$.
Consider the subcase when there exist $b_{j}\in T_{1}$, $j=2,3$ such that
$h(b_{2})\in(u_{2})$ and $h(b_{3})\in(f_{3})$ but
$h(T_{1})\cap(u_{3},u^{\prime}_{3})=\emptyset$. As above we may suppose that
after several procedures we changed $P_{b}$ such that $b_{j}\in(f_{j})$ and
the new $c^{\prime}_{2}$ is the old $h(b_{2})$. If $h(b_{2})\not\in C_{3}\cup
C_{2}$ then we may suppose that $h(b_{2})\not\in(w_{12})$. As in the first
case we may change $b_{3}$ such that $h(b_{3})\not\in(w_{23})$. Indeed, the
only problem could be if the old $h(b_{3})\in\\{u_{3},u^{\prime}_{3}\\}$,
which is not the case. We have no obstruction to change as usual $b_{1}$ such
that $h(b_{1})\not\in(w_{13})$ and so note that the interval
$[f_{2},h(b_{2})]$ (resp. $[f_{3},h(b_{3})]$, or $[f_{1},h(b_{1})]$) has at
most $w_{23}$ (resp. $w_{13}$, or $w_{12}$) from $W$. Thus the intervals
$[f_{j},h(b_{j})]$, $j\in[3]$ are disjoint.
If $h(b_{2})\in C_{3}$ then either $b_{2}=w_{23}$, or $w_{12}$. But
$b_{2}\not=w_{23}$ because otherwise $h(w_{23})\in(u_{2})$ contradicting
$(*)$. Similarly, $b_{2}\not=w_{12}$. If $h(b_{2})=w_{12}$ (resp.
$h(b_{2})=w_{23}$) then $b_{2}\not=w_{23}$ (resp. $b_{2}\not=w_{12}$) because
otherwise we get a contradiction with $(*)$. Thus $w_{12}$ (resp. $w_{23}$) is
the only monomial of $W$ which belongs to $[f_{2},h(b_{2})]$. Choosing $b_{3}$
such that $h(b_{3})\not\in(w_{23})$ (resp. $h(b_{3})\not\in(w_{12})$) and
$b_{1}$ such that $h(b_{1})\not\in(w_{12})$ (resp. $h(b_{1})\not\in(w_{13})$)
we get disjoint the corresponding intervals.
Now consider the subcase when there exist $b_{j}\in T_{1}$, $j=2,3$ such that
$h(b_{2})\in(u_{2})$ and $h(b_{3})\in(u_{3})$. If $h(b_{2})\not\in(f_{3})$ and
$h(b_{3})\not\in(f_{2})$ then as above we may assume that with a different
$P_{b}$, if necessary, we may reduce to the subcase when $b_{j}\in(f_{j})$,
$j=2,3$. In general this is not simple because $h(b_{2})$ as in Example 5 can
have no divisors from $B\cap(f_{2})$, which are not in
$\\{u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3}\\}$ and there exist no other
$c\in U_{1}$ multiple of $u_{2}$. In such situation we are force to remain on
the old $c^{\prime}_{2}$ taking
$T^{\prime}_{1}=T_{1}\cup\\{u_{2},u^{\prime}_{2}\\}$ and
$U^{\prime}_{1}=U_{1}\cup\\{c^{\prime}_{2}\\}$. If there exists a bad path
starting on a divisor from
$B\setminus\\{u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3}\\}$ of
$c^{\prime}_{2}$ then we go to case 3). Otherwise, we should consider also the
pathes starting with the divisors of $c^{\prime}_{2}$ from
$B\setminus\\{u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3}\\}$ completing
$T^{\prime}_{1}$ to ${\tilde{T}}_{1}$. Note that because of $2^{\prime})$ we
may speak now about $T_{m}$ instead $T_{1}$.
Changing in $P_{b}$ the intervals $[f_{j},c^{\prime}_{j}]$,
$[b_{j},h(b_{j})]$, $j=2,3$ with $[f_{j},h(b_{j})]$,
$[u^{\prime}_{j},c^{\prime}_{j}]$, $j=2,3$ we may assume the new
$c^{\prime}_{2},c^{\prime}_{3}$ are in $U_{m}=h(T_{m})$ for some $m$ and the
proof goes as above. If let us say $h(b_{2})\in(f_{3})$ then we must be
carefully since it is possible that the new intervals $[f_{j},c^{\prime}_{j}]$
could be not disjoint. A nice subcase is for example when $h(b_{2})$ is a
least common multiple of $u_{2},u_{3}$, which we study below.
If $w_{23}\in B$ then we we may suppose $u_{2}=w_{23}$. Indeed, if
$w_{23}\not\in\\{u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3}\\}$ and
$u_{2}=x_{p}f_{2}$ for some $p\not\in\operatorname{supp}w_{23}$ then
$h(b_{2})=x_{p}w_{23}$. Since $b_{2}\not=u_{2}$ and $b_{2}\in(f_{2})$ it
follows that $b_{2}=w_{23}$. But this contradicts the property $(*)$. Suppose
that $a_{t_{1}}=b_{2}$. Then note that $h(b_{2})=c_{t_{1}}=x_{p}w_{23}$ for
some $p$ and it follows that $a_{t_{1}}=b_{2}=x_{p}f_{2}$ since
$a_{t_{1}}\not=w_{23}=u_{2}$. Changing in $P_{b}$ the intervals
$[f_{2},c^{\prime}_{2}]$, $[b_{2},h(b_{2})]$ with $[f_{2},h(b_{2})]$,
$[u^{\prime}_{2},c^{\prime}_{2}]$, we may assume the new $c^{\prime}_{2}$ is
in $U_{m}$. We claim that $w_{23}$ is the only monomial from $B\cap W$ which
is in $[f_{2},c^{\prime}_{2}]$. Indeed, $w_{12}$ could be another monomial
from $B\cap W$ which is present in the new $[f_{2},c^{\prime}_{2}]$. This
could be true only if $a_{t_{1}}=w_{12}$. Thus $h(w_{12})=c_{t_{1}}\in(u_{2})$
which is not possible again by $(*)$. The same procedure we use to include a
new $c^{\prime}_{3}$ in $U_{m}$. Since $u_{2}=w_{23}$ cannot be among
$u_{3},u^{\prime}_{3}$ we see that only $w_{13}$ could be among them. Suppose
that $u_{3}=w_{13}$. Clearly the new $[f_{3},c^{\prime}_{3}]$ cannot contain
$w_{23}$. Choose as in the first case $b_{1}\in(f_{1})$ such that
$h(b_{1})\not\in(w_{13})$ and the new intervals $[f_{j},c^{\prime}_{j}]$,
$j\in[3]$ are disjoint. If $w_{13}\not\in\\{u_{3},u^{\prime}_{3}\\}$ then we
might have only $b_{3}=w_{13}$ and we may repeat the argument.
A problem could appear when the new $[f_{j},c^{\prime}_{j}]$, $j=2,3$ contain
$w_{12}$, $w_{13}$ because then we may not find $b_{1}$ as before. Note that
this problem could appear only when
$w_{12},w_{13}\in\\{u_{2},u^{\prime}_{2},u_{3},u^{\prime}_{3}\\}$ because of
$(*)$. We will change the new $c^{\prime}_{2}$ such that will not belong to
$(f_{1})$. Changing $P_{b}$ we may suppose that $b_{j}\in(f_{j})$, $j=2,3$
(again this change is not so simple as we saw above). We have
$h(b_{2})=c_{t_{1}}=x_{p}w_{12}$ for some $p$ and it follows that
$a_{t_{1}}=b_{2}=x_{p}f_{2}$ since $a_{t_{1}}\not=w_{12}=u_{2}$. Suppose that
$t_{1}>1$. Thus $a_{t_{1}}|c_{t_{1}-1}$ and we see that $c_{t_{1}-1}$ is not
in $(f_{1})$ because otherwise we get $c_{t_{1}-1}=x_{p}w_{12}=c_{t_{1}}$,
which is false. If $a_{t_{1}-1}\in(f_{2})$ then changing in $P_{b}$ the
intervals $[a_{t_{1}},c_{t_{1}}]$, $[a_{t_{1}-1},c_{t_{1}-1}]$,
$[f_{2},c^{\prime}_{2}]$ by $[f_{2},c_{t_{1}-1}]$, $[u_{2},c_{t_{1}}]$,
$[u^{\prime}_{2},c^{\prime}_{2}]$ we see that the new $c^{\prime}_{2}$ is not
in $(f_{1})$ and belongs to $U_{m}$. If $w_{12}\in C$ then we get
$h(b_{2})=w_{12}$ and the above argument works again, $c_{t_{1}-1}$ being the
new $c^{\prime}_{2}$.
When $a_{t_{1}-1}\not\in(f_{2})$ but $u^{\prime}_{2}|c_{t_{1}-1}$ we reduce
the problem to the subcase when the path $\\{a_{1},\ldots,a_{t_{1}-1}\\}$ goes
from $a_{1}$ to $u^{\prime}_{2}$ and now $u^{\prime}_{2}\not\in(f_{1})$. As
above we may change $P_{b}$ such that the new $b_{2}=a_{t_{1}-1}\in(f_{2})$
and the new $c^{\prime}_{2}$, that is the old $c_{t_{1}-1}$ is not in
$(f_{1})$.
If $a_{t_{1}-1}\not\in(f_{2})$, $u^{\prime}_{2}\not|c_{t_{1}-1}$ but there
exists ${\tilde{a}}\in B\cap(f_{2})$ a divisor of $c_{t_{1}-1}$ then
${\tilde{a}}\not=u_{2}$ because otherwise we get $c_{t_{1}-1}=c_{t_{1}}$. Now
we repeat the first part of the case 2). If ${\tilde{a}}=a_{v}$ for some
$1\leq v<t_{1}-1$ then changing in $P_{b}$ the intervals $[a_{p},c_{p}]$,
$v\leq p<t_{1}$ by $[a_{v},c_{t_{1}-1}]$, $[a_{p+1},c_{p}]$, $v\leq p<t_{1}-1$
we see that the new $c_{v}$ (resp. $c_{v+1}$) is the old $c_{t_{1}-1}$ (resp.
$c_{t_{1}}$)). Now changing the intervals $[a_{v},c_{v}]$,
$[a_{v+1},c_{v+1}]$, $[f_{2},c^{\prime}_{2}]$ by $[f_{2},c_{v}]$,
$[w_{12},c_{v+1}]$, $[u^{\prime}_{2},c^{\prime}_{2}]$ we see that the new
$c^{\prime}_{2}\not\in(f_{1})$ and belongs to $U_{m}$. If
${\tilde{a}}\not\in\\{a_{1},\ldots,a_{t_{1}}\\}$ then we are in one of the
above subcases $1^{\prime})$, $2^{\prime})$ solved already.
We may use this argument to change $c^{\prime}_{j}$, $j=2,3$ such that it is
not in $(f_{1})$ anymore, but as long as $h(b_{j})\not=c_{1}$, that is the
corresponding $t_{1}>1$. However, we may have $h(b_{j})=c_{1}$ only for one
$j>1$, because if for instance $h(b_{3})=c_{1}$ then $c_{1}\in C_{2}\cup
C_{3}$. If $c_{1}\in C_{3}$ then we see that $w_{23}\in B$ and $a_{1}=w_{23}$.
But this contradicts $(*)$ because $h(w_{23})\in(u_{2})$. If
$c^{\prime}_{2}\in C_{2}$ then $c^{\prime}_{2}=w_{23}$ and either
$a_{1}\in(f_{2})$, or $a_{1}\in(f_{3})$, that is $a_{1}$ cannot be $b_{2}$ and
$b_{3}$ in the same time. Thus at least one of the new $c^{\prime}_{j}$,
$j=2,3$ could be taken $\not\in(f_{1})$. If let us say only
$c^{\prime}_{3}\in(f_{1})$ then choose $b_{1}\in T_{1}\cap(f_{1})$ such that
$h(b_{1})\not\in(w_{13})$ as before. The interval $[f_{1},h(b_{1})]$ is
disjoint from the other new constructed intervals, which is enough as we saw
in case 1).
In the third case, let $a_{1},\ldots,a_{t_{1}}$ be a bad path starting with
$a_{1}$. Set $c_{j}=h(a_{j})$, $j\in[t_{1}]$. Then $c_{t_{1}}=bx_{l_{1}}$ and
let us say $b=f_{1}x_{i}$. If $a_{t_{1}}\in(f_{1})$ then changing in $P_{b}$
the interval $[a_{t_{1}},c_{t_{1}}]$ by $[f_{1},c_{t_{1}}]$ we get a partition
on $I/J$ with sdepth $d+2$. Thus we may assume that $a_{t_{1}}\not\in(f_{1})$.
If $f_{1}x_{l_{1}}\in\\{a_{1},\ldots,a_{t_{1}-1}\\}$, let us say
$fx_{l_{1}}=a_{v}$, $1\leq v<t_{1}$ then we may replace in $P_{b}$ the
intervals $[a_{p},c_{p}],v\leq p\leq t_{1}$ with the intervals
$[a_{v},c_{t_{1}}],[a_{p+1},c_{p}],v\leq p<t_{1}$. Now we see that we have in
$P_{b}$ the interval $[f_{1}x_{l_{1}},f_{1}x_{i}x_{l_{1}}]$ and switching it
with the interval $[f_{1},f_{1}x_{i}x_{l_{1}}]$ we get a partition with sdepth
$\geq d+2$ for $I/J$.
Thus we may assume that $f_{1}x_{l_{1}}\notin\\{a_{1},...,a_{t_{1}}\\}$. Now
set $a_{t_{1}+1}=fx_{l_{1}}$. Let
$a_{t_{1}+1},\ldots,a_{k}$ be a path starting with $a_{t_{1}+1}$ and set
$c_{j}=h(a_{j})$, $t_{1}<j\leq k$. If $a_{p}=a_{v}$ for $v\leq t_{1}$,
$p>t_{1}$ then change in $P_{b}$ the intervals $[a_{j},c_{j}],v\leq j\leq p$
with the intervals $[a_{v},c_{p}],[a_{j+1},c_{j}],v\leq j<p$. We have in
$P_{b}$ an interval $[f_{1}x_{l_{1}},f_{1}x_{i}x_{l_{1}}]$ and switching it to
$[f_{1},f_{1}x_{i}x_{l_{1}}]$ we get a partition with sdepth $\geq d+2$ for
$I/J$. Thus we may suppose that in fact
$a_{p}\not\in\\{b,a_{1},\ldots,a_{p-1}\\}$ for any $p>t_{1}$ (with respect to
any path starting with $a_{t_{1}+1}$). We have three subcases:
$1^{\prime\prime})$ there exist no weak path and no bad path starting with
$a_{t_{1}+1}$,
$2^{\prime\prime})$ there exists a weak path starting with $a_{t_{1}+1}$ but
no bad path starts with $a_{t_{1}+1}$,
$3^{\prime\prime})$ there exists a bad path starting with $a_{t_{1}+1}$.
Set $T_{2}=\\{b^{\prime}\in B:\mbox{there\ exists\ a\ path}\ \
a_{t_{1}+1},\ldots,a_{k}\ \ \mbox{with}\ \ a_{k}=b^{\prime}\\}$. We treat the
subcases $1^{\prime\prime})$, $2^{\prime\prime})$ as the cases 1), 2) and find
$I^{\prime}$ generated by a subset of $\\{f_{1},f_{2},f_{3}\\}\cup G_{2}$, or
$\\{f_{1},f_{2},f_{3}\\}\cup G^{\prime}_{2}$, or
$\\{f_{1},f_{2},f_{3}\\}\cup{\tilde{G}}_{2}$, where $G_{2}$, $G^{\prime}_{2}$,
${\tilde{G}}_{2}$, are obtained from $T_{2}$ and as above $T^{\prime}_{2}$, or
${\tilde{T}}_{2}$.
In the subcase $3^{\prime\prime})$, let $a_{t_{1}+1},\ldots,a_{t_{2}}$ be a
bad path starting with $a_{t_{1}+1}$ and set $c_{j}=h(a_{j})$ for $j>t_{1}$.
We saw that the whole set $\\{a_{1},\ldots,a_{t_{2}}\\}$ has different
monomials. As above $c_{t_{2}}=bx_{l_{2}}$ and we may reduce to the case when
$f_{1}x_{l_{2}}\not\in\\{a_{1},\ldots,a_{t_{1}}\\}$. Set
$a_{t_{2}+1}=f_{1}x_{l_{2}}$ and again we consider three subcases, which we
treat as above. Anyway after several such steps we must arrive in the case
$p=t_{m}$ when either we may proceed as in the subcases $1^{\prime\prime})$,
$2^{\prime\prime})$, or $b|c_{t_{m}}$ and again a certain $f_{1}x_{l_{m}}$ is
not among $\\{a_{1},\ldots,a_{t_{m}}\\}$ and taking
$a_{t_{m}+1}=f_{1}x_{l_{m}}$ there exist no bad path starting with
$a_{t_{m}+1}$. This follows since we may reduce to the subcase when the set
$\\{a_{1},\ldots,a_{t_{m}}\\}$ has different monomials and so the procedures
should stop for some $m$. Finally, using
$T_{m}=\\{b^{\prime}\in B:\mbox{there\ exists\ a\ path}\ \
a_{t_{m}+1},\ldots,a_{k}\ \ \mbox{with}\ \ a_{k}=b^{\prime}\\}$
(resp. $T^{\prime}_{m}$, or ${\tilde{T}}_{m}$) as $T_{1}$ (resp.
$T^{\prime}_{1}$, or ${\tilde{T}}_{1}$) above we are done.
Proof of Theorem 6. By Theorems 2, 3 we may suppose that $2r\leq s\leq q+r$
and we may assume that $E$ contains only monomials of degrees $d+1$ by [15,
Lemma 1.6]. Apply induction on $|E|$. If $E=\emptyset$ we may apply Theorem 4.
Suppose that $|E|>0$ and $B\cap(f_{1},\ldots,f_{r})\not=\emptyset$, $r=2,3$,
otherwise we get $\operatorname{depth}_{S}I/J\leq d+1$ using [15, Lemma 1.5]
applied to any $f_{i}$. We may choose $b\in B\cap(f_{1},f_{2},f_{3})$ which is
not in $W$ if $r=2,3$ and $|B\cap(f_{1},\ldots,f_{r})|>3\geq|B\cap W|$.
However, $|B\cap(f_{1},\ldots,f_{r})|<2r$ gives
$\operatorname{depth}_{S}(f_{1},\ldots,f_{r})/J\cap(f_{1},\ldots,f_{r})\leq
d+1$ by Theorem 3 and it follows that $\operatorname{depth}_{S}I/J\leq d+1$
using the Depth Lemma applied to the exact sequence
$0\rightarrow(f_{1},\ldots,f_{r})/J\cap(f_{1},\ldots,f_{r})\rightarrow
I/J\rightarrow(E)/(J,f_{1},\ldots,f_{r})\cap(E)\rightarrow 0.$
Thus if $r=2,3$ we may suppose to find $b\in
B\cap(f_{1},\ldots,f_{r})\setminus W$. Renumbering $f_{i}$ we may suppose that
$b\in(f_{1})\setminus(f_{2},\ldots,f_{3})$.
Apply induction on $r\leq 3$. Using Theorem 4 and induction hypothesis on
$|E|$ and $r$ apply [17, Lemma 1.1]. Thus we may suppose that
$C\subset((f_{1})\cap(f_{2}))\cup((E)\cap(f_{1},f_{2}))\cup(\cup_{a,a^{\prime}\in
E,a\not=a^{\prime}}(a)\cap(a^{\prime}))$, if $r=2$, or
$C\subset(\cup_{i,j\in[3],i\not=j}(f_{i})\cap(f_{j}))\cup((E)\cap(f_{1},f_{2},f_{3}))\cup(\cup_{a,a^{\prime}\in
E,a\not=a^{\prime}}(a)\cap(a^{\prime}))$ if $r=3$.
Set $I^{\prime}_{b}=(f_{2},\ldots,f_{r},B\setminus\\{b\\})$,
$J^{\prime}_{b}=I^{\prime}_{b}\cap J$. Clearly $b\not\in I^{\prime}_{b}$ and
so in the following exact sequence
$0\rightarrow I^{\prime}_{b}/J^{\prime}_{b}\rightarrow I/J\rightarrow
I/(J,I^{\prime}_{b})\rightarrow 0$
the last term has depth $\geq d+1$. If the first term has sdepth $\leq d+1$
then it has depth $\leq d+1$ by induction hypothesis on $r$, case $r=1$ being
done in Theorem 4. Thus we may suppose that
$\operatorname{sdepth}_{S}I^{\prime}_{b}/J^{\prime}_{b}\geq d+2$ and we may
apply Lemmas 5, 6. Then we get either $\operatorname{sdepth}_{S}I/J\geq d+2$
contradicting our assumption, or there exists a nonzero ideal
$I^{\prime}\subsetneq I$ generated by a subset $G$ of $B$, or by $G$ and a
subset of $\\{f_{1},f_{2},f_{3}\\}$ such that
$\operatorname{sdepth}_{S}I^{\prime}/J^{\prime}\leq d+1$ for $J^{\prime}=J\cap
I^{\prime}$ and $\operatorname{depth}_{S}I/(J,I^{\prime})\geq d+1$. In the
last case we see that $\operatorname{depth}_{S}I^{\prime}/J^{\prime}\leq d+1$
by induction hypothesis on $r$, $|E|$, or by Theorem 4 and so
$\operatorname{depth}_{S}I/J\leq d+1$ by the Depth Lemma applied to the
following exact sequence
$0\rightarrow I^{\prime}/J^{\prime}\rightarrow I/J\rightarrow
I/(J,I^{\prime})\rightarrow 0.\ \ \ \ \ \ \ \Box$
The following bad example it is useful to illustrate somehow our proof.
###### Example 5.
Let $n=6$, $r=3$, $d=1$, $f_{i}=x_{i}$ for $i\in[3]$,
$E=\\{x_{4}x_{5},x_{5}x_{6}\\}$, $I=(x_{1},x_{2},x_{3},E)$ and
$J=(x_{2}x_{4},x_{3}x_{4},x_{1}x_{2}x_{6},x_{1}x_{3}x_{6},x_{1}x_{4}x_{6},x_{1}x_{5}x_{6},x_{2}x_{3}x_{6},x_{2}x_{5}x_{6},$
$x_{3}x_{5}x_{6})$. Then
$B=\\{x_{1}x_{2},x_{1}x_{3},x_{2}x_{3},x_{1}x_{4},x_{1}x_{5},x_{2}x_{5},x_{3}x_{5}\\}\cup
E$ and
$C=\\{x_{1}x_{2}x_{3},x_{1}x_{2}x_{5},x_{2}x_{3}x_{5},x_{1}x_{3}x_{5},x_{1}x_{4}x_{5},x_{4}x_{5}x_{6}\\}.$
Take $b=x_{1}x_{4}$ and
$I_{b}=(x_{2},x_{3},x_{1}x_{2},x_{1}x_{3},x_{1}x_{5},E)$, $J_{b}=I_{b}\cap J$.
There exists a partition $P_{b}$ with sdepth $3$ on $I_{b}/J_{b}$ given by the
intervals $[x_{2},x_{1}x_{2}x_{3}]$, $[x_{3},x_{1}x_{3}x_{5}]$,
$[x_{1}x_{5},x_{1}x_{2}x_{5}]$, $[x_{2}x_{5},x_{2}x_{3}x_{5}]$,
$[x_{4}x_{5},x_{1}x_{4}x_{5}]$, $[x_{5}x_{6},x_{4}x_{5}x_{6}]$. We have
$c^{\prime}_{2}=x_{1}x_{2}x_{3}$, $c^{\prime}_{3}=x_{1}x_{3}x_{5}$ and
$u_{2}=x_{2}x_{3}$, $u^{\prime}_{2}=x_{1}x_{2}$, $u_{3}=x_{3}x_{5}$,
$u^{\prime}_{3}=x_{1}x_{3}$. Clearly, $u_{2}=w_{23}$.
Take $a_{1}=x_{1}x_{5}$, $c_{1}=x_{1}x_{2}x_{5}$, $a_{2}=x_{2}x_{5}$,
$c_{2}=x_{2}x_{3}x_{5}$. The path $\\{a_{1},a_{2}\\}$ is maximal weak because
the divisors from $B$ of $c_{2}$ are $a_{2},u_{2},u_{3}$. Then
$T_{1}=\\{a_{1},a_{2}\\}$ and we change in $P_{b}$ as in the proof the
intervals $[x_{2},c^{\prime}_{2}]$, $[a_{2},c_{2}]$ by $[x_{2},c_{2}]$,
$[u^{\prime}_{2},c^{\prime}_{2}]$. Thus the new $c^{\prime}_{2}$ is the old
$c_{2}$. Now note that this new $c^{\prime}_{2}$ is a multiple of $u_{3}$ and
it is the only monomial from $h(T_{1})$, which is a such multiple. Thus we had
to take $u_{3}$ in the new $T^{\prime}_{1}$, and $u^{\prime}_{3}$ as well and
certainly $c^{\prime}_{3}$ is added to $h(T_{1})$. Clearly, all divisors from
$B$ of $c^{\prime}_{3}$ are in
$T^{\prime}_{1}=T_{1}\cup\\{u_{3},u^{\prime}_{3}\\}$. But the former
$u^{\prime}_{2}$ divides $c_{1}$ and so should be added to $T^{\prime}_{1}$.
Thus we have $I^{\prime}=(b,E)$, $J^{\prime}=J\cap I^{\prime}$ and
$I/(J,I^{\prime})$ has a partition of sdepth $3$ given by the intervals
$[x_{2},c_{2}]$, $[x_{3},c^{\prime}_{3}]$, $[x_{1},x_{1}x_{2}x_{5}]$. If
$\operatorname{sdepth}_{S}I^{\prime}/J^{\prime}\geq 3$ then we get
$\operatorname{sdepth}_{S}I/J\geq 3$, which is false. Otherwise,
$\operatorname{sdepth}_{S}I^{\prime}/J^{\prime}\leq 2$ and we get
$\operatorname{depth}_{S}I^{\prime}/J^{\prime}\leq 2$ by [12, Theorem 4.3] and
so $\operatorname{depth}_{S}I/J\leq 2$ using the Depth Lemma.
## References
* [1] W. Bruns and J. Herzog, Cohen-Macaulay rings, Revised edition. Cambridge University Press (1998).
* [2] W. Bruns, C. Krattenthaler, J. Uliczka, Stanley decompositions and Hilbert depth in the Koszul complex, J. Commutative Alg., 2 (2010), 327-357.
* [3] W. Bruns, J. J. Moyano-Fernández, J. Uliczka, Hilbert regularity of ZZ-graded modules over polynomial rings, arXiv:AC/1308.2917, 2013.
* [4] M. Cimpoeas, The Stanley conjecture on monomial almost complete intersection ideals, Bull. Math. Soc. Sci. Math. Roumanie, 55(103) (2012), 35-39.
* [5] J. Herzog, M. Vladoiu, X. Zheng, How to compute the Stanley depth of a monomial ideal, J. Algebra, 322 (2009), 3151-3169.
* [6] B. Ichim, J. J. Moyano-Fernández, How to compute the multigraded Hilbert depth of a module, to appear in Mathematische Nachrichten, arXiv:AC/1209.0084.
* [7] B. Ichim, A. Zarojanu, An algorithm for computing the multigraded Hilbert depth of a module, arXiv:AC/1304.7215.
* [8] A. Popescu, Special Stanley Decompositions, Bull. Math. Soc. Sc. Math. Roumanie, 53(101), no 4 (2010), 363-372, arXiv:AC/1008.3680.
* [9] A. Popescu, An algorithm to compute the Hilbert depth , J. Symb. Comput.,66, (2015), 1-7, arXiv:AC/1307.6084.
* [10] D. Popescu, Stanley conjecture on intersections of four monomial prime ideals, Communications in Alg., 41 (2013), 4351-4362, arXiv:AC/1009.5646.
* [11] D. Popescu, Graph and depth of a square free monomial ideal, Proceedings of AMS,140 (2012), 3813-3822, arXiv:1104.5596v1, 2011.
* [12] D. Popescu, Depth of factors of square free monomial ideals, to appear in Proceedings AMS, arXiv:AC/1110.1963.
* [13] D. Popescu, Upper bounds of depth of monomial ideals, J. Commutative Algebra, 5, 2013, 323-327, arXiv:AC/1206.3977.
* [14] D. Popescu, Depth of a monomial ideal, in Proceedings of the Seventh Congress of Romanian Mathematicians, 2011.
* [15] D. Popescu, A. Zarojanu, Depth of some square free monomial ideals, Bull. Math. Soc. Sci. Math. Roumanie, 56(104) (2013),117-124.
* [16] D. Popescu, A. Zarojanu, Depth of some special monomial ideals, Bull. Math. Soc. Sci. Math. Roumanie, 56(104) (2013), 365-368, arxiv:AC/1301.5171v1.
* [17] D. Popescu, A. Zarojanu, Three generated, squarefree, monomial ideals , to appear in Bull. Math. Soc. Sci. Math. Roumanie, 58(106), no 3, (2015), arXiv:AC/1307.8292v6
* [18] A. Rauf, Depth and Stanley depth of multigraded modules, Comm. Algebra, 38 (2010),773-784.
* [19] Y.H. Shen, Lexsegment ideals of Hilbert depth 1 (2012), arXiv:AC/1208.1822v1.
* [20] R. P. Stanley, Linear Diophantine equations and local cohomology, Invent. Math. 68 (1982) 175-193.
* [21] J. Uliczka, Remarks on Hilbert series of graded modules over polynomial rings, Manuscripta Math., 132 (2010), 159-168.
* [22] A. Zarojanu, Stanley Conjecture on intersection of three monomial primary ideals, Bull. Math. Soc. Sc. Math. Roumanie, 55(103),(2012), 335-338.
|
arxiv-papers
| 2013-09-19T14:10:31 |
2024-09-04T02:49:51.158408
|
{
"license": "Public Domain",
"authors": "Adrian Popescu and Dorin Popescu",
"submitter": "Adrian Popescu",
"url": "https://arxiv.org/abs/1309.4986"
}
|
1309.5000
|
# Optimally focused cold atom systems obtained using density-density
correlations
Andika Putra Joint Quantum Institute, University of Maryland, and National
Institute of Standards and Technology, College Park, Maryland, 20742, USA
Daniel L. Campbell Joint Quantum Institute, University of Maryland, and
National Institute of Standards and Technology, College Park, Maryland, 20742,
USA Ryan M. Price Joint Quantum Institute, University of Maryland, and
National Institute of Standards and Technology, College Park, Maryland, 20742,
USA Subhadeep De Joint Quantum Institute, University of Maryland, and
National Institute of Standards and Technology, College Park, Maryland, 20742,
USA CSIR-National Physical Laboratory, New Delhi 110012, India I. B.
Spielman Joint Quantum Institute, University of Maryland, and National
Institute of Standards and Technology, College Park, Maryland, 20742, USA
###### Abstract
Resonant absorption imaging is a common technique for detecting the two-
dimensional column density of ultracold atom systems. In many cases, the
system’s thickness along the imaging direction greatly exceeds the imaging
system’s depth of field, making the identification of the optimally focused
configuration difficult. Here we describe a systematic technique for bringing
Bose-Einstein condensates (BEC) and other cold-atom systems into an optimal
focus even when the ratio of the thickness to the depth of field is large: a
factor of 8 in this demonstration with a BEC. This technique relies on
defocus-induced artifacts in the Fourier-transformed density-density
correlation function (the power spectral density, PSD). The spatial frequency
at which these artifacts first appear in the PSD is maximized on focus; the
focusing process therefore both identifies and maximizes the range of spatial
frequencies over which the PSD is uncontaminated by finite-thickness effects.
## I Introduction
Since the most important technique for obtaining properties of ultracold atoms
is direct imaging, a well-designed and well-aligned imaging system is crucial
for obtaining high quality data which is valid at all length scales. While
large scale properties such as the system’s width or peak density can be
obtained with little effort, significant care must be taken for experiments
requiring very good spatial resolution Bakr et al. (2010); Sherson et al.
(2010), or those studying correlations Folling et al. (2005); Hung et al.
(2011). It is difficult to bring objects extended along the imaging axis, such
as degenerate Fermi gases Granade et al. (2002); Regal et al. (2003), 3D Mott
insulators Greiner et al. (2002), and Bose-Einstein condensates (BECs)
Anderson et al. (1995); Davis et al. (1995), into focus particularly after
time-of-flight (TOF) expansion because their spatial thickness often exceeds
the imaging system’s depth of field. Even for such objects, a high degree of
accuracy in focusing is required to minimize imaging artifacts. Understanding
and minimizing these artifacts is particularly important when studying
density-density correlations, where the artifacts can be confused with the
correlation signal under study Hung et al. (2011); Choi et al. (2012); Seo et
al. (2013); Langen (2013); De et al. (2012). Here we describe a fairly generic
technique for focusing on these extended objects which is far more precise
than simply optimizing the “sharpness” of imaged atom clouds.
Absorption imaging is a ubiquitous approach for measuring the density
distribution of ultracold atom systems Ketterle et al. (1999). A probe beam
illuminates the atomic system, and the resulting shadow is imaged onto a
scientific camera, typically a charge-coupled device (CCD) or complementary
metal-oxide-semiconductor (CMOS) detector. Ideally, the fraction of light
absorbed would be directly related to the two-dimensional column density
$\rho_{\rm 2D}(x,y)=\int{\rm d}z\;\rho(x,y,z)$ of the atoms along the imaging
direction ${\bf e}_{z}$, where $\rho(x,y,z)$ is the density of atoms. If the
thickness $\delta z$ along ${\bf e}_{z}$ exceeds the imaging system’s depth of
field, then some of the atomic distribution must necessarily be out of focus,
invalidating any simple relationship between absorption and column density.
Given this, it is a challenge to obtain the optimal focal plane of the
extended system that minimizes the artifacts resulting from this defocus,
e.g., at the center of a distribution symmetric along ${\bf e}_{z}$.
Typically a system is brought into a focus by minimizing the size or apparent
diffraction effects from a compact object such as a trapped BEC; in many
cases, no such compact reference at the desired image plane is available. In
this paper, we present a technique for determining the optimal focus of
absorption-imaged extended objects. Using this technique, we identify the
focal plane within an accuracy of $2\,\mathrm{\mu m}$ for a $\delta
z=150\,\mathrm{\mu m}$ thick object. Specifically, given an object with
density-density correlations Hung et al. (2011) with a spatial correlation
length $\ell$, we show that observations of correlations in the optical
absorption as a function of camera position allow us to bring the object into
focus to within a fraction of the depth of field associated with $\ell$, even
without knowing the details of the correlation function. This optimal focus is
the camera position where the imaged auto-correlation function (ACF) most
accurately reflects the atomic density-density correlations, minimizing both
defocus-induced artifacts and the resolution limiting effect of the system’s
finite thickness Langen (2013); De et al. (2012).
In this paper, we review the basic theoretical formulation required to
understand light propagating through an absorbing dielectric medium. We then
consider several example images created by different idealized objects, in
each case noting how to determine their optimal focus. Lastly, we
experimentally apply this technique to images of BECs after TOF.
## II Theory
Monochromatic light of free-space wavelength $\lambda$ and wavenumber
$k_{0}=2\pi/\lambda$ propagating through an object with complex relative
permittivity $\varepsilon(\mathbf{r})=\epsilon/\epsilon_{0}$ and relative
susceptibility $\chi({\bf r})=\varepsilon(\mathbf{r})-1$, where $\epsilon$ is
the permittivity, and $\epsilon_{0}$ is the electric constant, is described by
the vectorial wave equation for the electric field $\mathbf{E(r)}$:
$\nabla^{2}\mathbf{E(r)}+{k_{0}^{2}}\varepsilon(\mathbf{r})\mathbf{E(r)}=-\nabla\left[\mathbf{E(r)}\cdot\nabla\ln\varepsilon(\mathbf{r})\right].$
(1)
In a medium where $\varepsilon(\mathbf{r})$ is slowly varying, the right-hand
side (rhs) of Eq. (1) can be neglected, reducing Eq. (1) to separate scalar
wave equations
$\nabla^{2}E\mathbf{(r)}+k_{0}^{2}\varepsilon(\mathbf{r})E\mathbf{(r)}=0,$
for each vector component of $\mathbf{E(r)}$, e.g., we might have
$\mathbf{E(r)}=E\mathbf{(r)}\mathbf{e}_{x}$ for light linearly polarized along
$\mathbf{e}_{x}$.
### II.1 Wavefield propagation
Here, we cast the above scalar wave equation into the form
$\displaystyle-\frac{\partial^{2}E\mathbf{(r)}}{\partial z^{2}}$
$\displaystyle=$
$\displaystyle\left[\nabla^{2}_{\perp}+k_{0}^{2}\right]E\mathbf{(r)}+k_{0}^{2}\chi({\bf
r})E(\mathbf{r}),$ (2)
suitable for light predominantly traveling along ${\bf e}_{z}$. For a known
field configuration at $E({\bf r})$ (such as the probe laser before it
interacts with the atoms), Eq. (2) has the formal solution
$E\left({\bf r}+\Delta z{\bf e}_{z}\right)=\exp\left[\pm i\Delta
z\sqrt{\nabla^{2}_{\perp}+k_{0}^{2}+k_{0}^{2}\chi({\bf
r})}\right]\\!E\left({\bf r}\right),$ (3)
describing the field propagated a distance $\Delta z$ along ${\bf e}_{z}$.
[plus minus sign]
Wave propagation in free space [i.e., $\chi({\bf r})=0$ in Eq.(2)] is solved
exactly in the angular spectrum representation Novotny and Hecht (2006)
$\displaystyle E_{\rm fs}({\bf r}\\!+\\!\Delta z{\bf e}_{z})$
$\displaystyle\\!\\!=\\!\\!$ $\displaystyle\mathbf{P}(\Delta z)E({\bf r})$
$\displaystyle\\!\\!=\\!\\!\\!$
$\displaystyle\int\\!\mathrm{d}^{2}\mathbf{k}_{\rm 2D}\\!\left[\tilde{{\bf
P}}(\mathbf{k}_{\rm 2D},\Delta z)\tilde{E}(\mathbf{k}_{\rm
2D},z)\right]\\!e^{i{\bf k}_{\rm 2D}\cdot\mathbf{r}_{\rm 2D}},$
for a forward going wave, with the 2D position ${\bf r}_{\rm
2D}=\left(x,y\right)$ and wavevector $\mathbf{k}_{\rm
2D}=\left(k_{x},k_{y}\right)$; the Fourier-transformed wavefield
$\tilde{E}({\bf k}_{\rm 2D},z)=\int{\rm d}^{2}{\bf r}_{\rm 2D}\>\exp(-i{\bf
k}_{\rm 2D}\cdot\mathbf{r}_{\rm 2D})E(\mathbf{r})$; and the transfer function
for propagating a distance $\Delta z$ in free space
$\tilde{{\bf P}}\left({\bf k}_{\rm 2D},\Delta z\right)=\exp\left[i\Delta
z\left(k_{0}^{2}-k^{2}_{\rm 2D}\right)^{1/2}\right].$
The transfer function behaves differently in two regions of spatial
frequencies: for $k^{2}_{\rm 2D}<k_{0}^{2}$, $\tilde{\bf P}$ is oscillatory
(propagating regime), and for $k^{2}_{\rm 2D}>k_{0}^{2}$, it is exponentially
decaying (evanescent regime).
Meanwhile, considering only $\chi({\bf r})$ [neglecting the first term in the
rhs of Eq. (2)], the absorption and refraction of light traveling a distance
$\Delta z$ is described by
$\displaystyle E_{\rm BL}({\bf r}\\!+\\!\Delta z{\bf e}_{z})$ $\displaystyle=$
$\displaystyle\mathbf{Q}(\Delta z)E({\bf r})$ $\displaystyle=$
$\displaystyle\exp\left[ik_{0}\int_{z}^{z+\Delta z}\mathrm{d}z\sqrt{\chi({\bf
r})}\right]E({\bf r}).$
Unlike the usual Beer-Lambert (BL, discussed in Sect. II.2), this expression
alone does not reflect a good approximation to beam propagation for systems of
any significant thickness.
### II.2 Beer-Lambert law and the paraxial approximation
To better understand the independent influence of the beam’s propagation and
its interaction with matter, we apply the paraxial approximation to Eq. (2),
allowing us to draw an analogue between the paraxial wave equation and the
Schrödinger equation, which can be solved numerically using a split-step
Fourier method (SSFM) J. A. Fleck, Jr. (1995). To understand the difference
between Eq. (II.1) and the usual BL law, we again turn to Eq. (2), now
assuming that the electric field can be written as $E({\bf
r})=\exp\left(ik_{0}z\right)E^{\prime}({\bf r})$, where $E^{\prime}({\bf r})$
is a slowly varying envelope along ${\bf e}_{z}$. Inserting this form into Eq.
(2) gives the paraxial wave equation
$\displaystyle-2ik_{0}\frac{\partial E^{\prime}\mathbf{(r)}}{\partial z}$
$\displaystyle=$
$\displaystyle\nabla^{2}_{\perp}E^{\prime}\mathbf{(r)}+k_{0}^{2}\chi({\bf
r})E^{\prime}(\mathbf{r}),$ (6)
where the assumed weak $z$ dependence of $E^{\prime}({\bf r})$ allowed us to
drop $\partial^{2}_{z}E^{\prime}({\bf r})$. Like above, the spatial evolution
of an initial $E^{\prime}({\bf r})$ can be partitioned into a spectral part
$\tilde{\bf P}^{\prime}({\bf k}_{\rm 2D},\Delta z)$ and a coordinate part
${\bf Q}^{\prime}(\Delta z)$, with
$\displaystyle\tilde{\bf P}^{\prime}({\bf k}_{\rm 2D},\Delta z)$
$\displaystyle=\exp\left(-i\frac{k_{\rm 2D}^{2}}{2k_{0}}\Delta z\right)$ (7)
$\displaystyle{\bf Q}^{\prime}(\Delta z)$
$\displaystyle=\exp\left[i\frac{k_{0}}{2}\int_{z}^{z+\Delta z}\chi({\bf
r})\mathrm{d}z\right].$ (8)
For the paraxial approximation to be valid, the condition $\left|\chi({\bf
r})\right|\ll 1$ must also hold: otherwise the ${\bf Q}^{\prime}(\Delta z)$
evolution would lead $E^{\prime}({\bf r})$ to depend strongly on $z$.
We numerically evolve the paraxial wave equation [Eq. (6)] along ${\bf e}_{z}$
using a split-step Fourier method (SSFM) Korpel et al. (1986); Feit and J. A.
Fleck (1978), where the operators in the rhs of Eq. (6) are split into two:
one operator represents wave propagation in a uniform medium using Eq. (7) and
the other operator takes into account the effect of refractive index variation
using Eq. (8). In the SSFM, we alternately apply the two evolution operators
with steps of size $\Delta z$. For each step, the complex amplitude
$E^{\prime}(\mathbf{r})$ is propagated first by $\mathbf{P}^{\prime}(\Delta
z/2)$, then by $\mathbf{Q}^{\prime}(\Delta z)$, and then again by
$\mathbf{P}^{\prime}(\Delta z/2)$. The resulting symmetrized split evolution
$E^{\prime}\left({\bf r}+\Delta z{\bf
e}_{z}\right)=\mathbf{P^{\prime}}\left(\Delta
z/2\right)\mathbf{Q^{\prime}}(\Delta z)\mathbf{P^{\prime}}\left(\Delta
z/2\right)E^{\prime}\left({\bf r}\right),$
has its first correction at order $\Delta z^{3}$.
The paraxial equations allow us to introduce the depth of field
$d_{\rm dof}=\frac{2k_{0}}{k_{\rm max}^{2}}=\frac{l_{\rm
min}^{2}}{\pi\lambda},$ (9)
where $k_{\rm max}^{2}$ is the largest $k_{\rm 2D}$ of interest and $l_{\rm
min}=2\pi/k_{\rm max}$ is the corresponding minimum length scale [these might
be specified by: the maximum significant wavevector in $\chi({\bf k}_{\rm
2D},z)$; the resolution of the physical imaging system; or at most by
$k_{0}$].
We obtain the BL law by assuming that the system is thin along ${\bf e}_{z}$,
i.e., both $\delta z\ll d_{\rm dof}$, and $\tilde{\bf P}^{\prime}({\bf k}_{\rm
2D},\Delta z)$ may be neglected. For purely absorbing materials where
$\chi({\bf r})\propto i\sigma_{0}\rho({\bf r})$, this gives the usual BL law
$\displaystyle I({\bf r}+\delta z{\bf e}_{z})$ $\displaystyle=$
$\displaystyle\exp\left[-\sigma_{0}\int_{z}^{z+\delta z}\rho({\bf
r})\mathrm{d}z\right]I({\bf r})$ (10)
describing the attenuation of the free space optical intensity $I({\bf
r})=c\epsilon_{0}\left|E({\bf r})\right|^{2}/2$ by absorbers of density
$\rho({\bf r})$ and scattering cross-section $\sigma_{0}$. This BL result can
also be obtained without the paraxial approximation by first neglecting the
$\nabla^{2}_{\perp}$ term in Eq. (2) (valid when $k_{\rm max}\delta z\ll 1$: a
more strict requirement than in the paraxial approximation where we had
$\delta z\ll d_{\rm dof}$) and again assuming $\left|\chi({\bf r})\right|\ll
1$, a small relative susceptibility 111Moreover, the gradient term
$\nabla\ln\varepsilon(\mathbf{r})$ in Eq. (1) cannot be safely neglected for
systems where $\left|\chi({\bf r})\right|$ is large or sufficiently rapidly
varying, although this would generally imply a breakdown of the paraxial
approximation as well..
In experiment, the BL law is generally applied by comparing the intensities
$I({\bf r}_{\rm 2D})$ and $I_{0}({\bf r}_{\rm 2D})$ measured with and without
atoms present, respectively. This relates the optical depth
${\rm OD}({\bf r}_{\rm 2D})\equiv-\ln\frac{I({\bf r}_{\rm 2D})}{I_{0}({\bf
r}_{\rm 2D})}=\sigma_{0}\rho_{\rm 2D}({\bf r}_{2D})$
to the 2D column density. In cold atom experiments, this column density is the
primary observable in experiment.
### II.3 Absorption Imaging
Figure 1: Dependence of BEC images on image plane position. (a) Intensity
corrected optical depth measured $d=-54\,\mathrm{\mu}{\rm m}$, 0 $\mu{\rm m}$,
and $54\,\mathrm{\mu}{\rm m}$ from the optimal focus: images (right) and line
cuts at $y=0$ (left). (b) The peak optical depth depends only weakly on $d$;
is not maximized at $d=0$; and has no structure on the $2\,\mathrm{\mu m}$
scale.
Here we consider systems of ultracold atoms illuminated by laser light on a
cycling transition, where the atom-light interaction is described by a complex
relative susceptibility
$\chi({\bf
r})=\frac{\sigma_{0}}{k_{0}}\left[\frac{i-2\delta/\Gamma}{1+{I}/{I_{\rm
sat}}+\left({2\delta}/{\Gamma}\right)^{2}}\right]\rho(\mathbf{r}).$
$\rho(\mathbf{r})$ is the atomic density; $\delta$ is the laser’s detuning
from atomic resonance; $\sigma_{0}=6\pi/k_{0}^{2}$ is the resonant scattering
cross-section; $\Gamma$ is the atomic linewidth; and $I_{\rm sat}$ is the
saturation intensity Ketterle et al. (1999). The standard BL law is valid for
dilute ($\rho\ll k_{0}^{3}$, see Ref. 222A dilute regime, with density
$\rho\ll k_{0}^{3}$, can be achieved by letting the atomic cloud to
ballistically expand.), spatially thin systems ($k_{0}\delta z\ll 1$),
illuminated by low intensity ($I_{0}\ll I_{\rm sat}$) probe beams.
The $I_{0}\ll I_{\rm sat}$ requirement can be lifted by introducing the
intensity-corrected optical depth
${\rm{OD_{cor}}}({\bf r}_{2D})\equiv-\ln\frac{I({\bf r}_{2D})}{I_{0}({\bf
r}_{2D})}+\frac{I_{0}({\bf r}_{2D})-I({\bf r}_{2D})}{I_{\rm sat}},$ (11)
which is related to the column density
$\rho_{\rm 2D}({\bf r}_{2D})=\frac{{\rm{OD_{cor}}}({\bf
r}_{2D})}{\sigma_{0}},$ (12)
of dilute ($\rho\ll k_{0}^{3}$), spatially thin systems ($k_{0}\delta z\ll
1$). Due to the limited dynamic range of the camera’s pixels Reinaudi et al.
(2007) and the presence of background light, it is technically difficult to
reliably detect uncorrected optical depths, larger than $\approx 4$. Thus, we
deliberately select $I_{0}>I_{\rm sat}$, saturating the transition with
$I_{0}$ such that ${\rm OD_{cor}}<3$.
In addition, the spatial thickness of many cold atom systems exceed the depth
of field leaving parts of its distribution along imaging direction inevitably
out of focus, thereby invalidating Eq. (12). Even for dilute clouds (after
sufficient TOF), images taken an equal distance above and below the focal
plane can differ. This lack of symmetry makes a straightforward determination
of the optimal focus difficult (lensing effects from even slightly off-
resonance imaging beams and aberrations in the imaging system can complicate
the situation further.)
To illustrate this difficulty, we consider images of BECs with the focal plane
displaced a distance $d=-54\,\mathrm{\mu m}$, 0 ${\rm\mu m}$, and
$54\,\mathrm{\mu m}$ from the BECs’ center (see Fig. 1). Because the BEC is
thick compared to the depth of field, Eq. (12) does not hold; in addition
lensing effects cause the cloud’s peak $\rm{OD_{cor}}$ to behave
asymmetrically when the focus is behind or in front of the cloud. In these
images, there are no sharp features that identify the optimal focus at the
micron level. Owing to the weak dependence of large-scale parameters such as
peak-height or width on defocus, such precise focusing is not required in many
experiments. As we see below, experiments that study correlations within such
images are extremely sensitive to defocus and new methods are required. Our
technique brings images such as these into focus, identifying an optimal focal
plane at the $\approx 2\,\mathrm{\mu m}$ level.
### II.4 Modeling
To obtain a basic understanding of our approach, we first consider the
defocused image of a $1\,\mathrm{\mu m}$ thick absorbing medium, inhomogenous
${\bf e}_{x}$-${\bf e}_{y}$ plane, bounded above and below by vacuum, with,
$\chi({\bf r})=ig(x,y)$ for $z\in(-0.5\,\mathrm{\mu m},0.5\,\mathrm{\mu m})$,
where $g(x,y)\geq 0$ is a Poisson distributed random variable. Like atoms
illuminated on resonance, this medium has a purely imaginary susceptibility.
The illuminating light is modeled by a plane wave with wavelength
$\lambda=780\,\mathrm{nm}$ suitable for imaging our 87Rb Bose-Einstein
condensates 333In our SSFM simulation of light traversing this medium, we used
a $\Delta z=1\,\mathrm{\mu m}$ step size.. While this object has no visible
structure, by virtue of its spectrally flat density-density correlation
function, it can be brought into focus.
The imaged intensity pattern $I(x,y)$ from this $1\,\mathrm{\mu m}$ layer
appears random at various distances from focus, but its correlations become
oscillatory. To reveal this information, we turn to its spatial power spectral
density: the magnitude squared of $I(x,y)$’s Fourier transform 444The Wiener-
Khincin theorem states that the spectral decomposition of the autocorrelation
function is equal to the power spectral density.. The PSD is circularly
symmetric in the spatial frequency $\mathbf{k}_{\rm 2D}=(k_{x},k_{y})$ plane.
Fig. 2a shows the PSD in this $k=\left|\mathbf{k}_{\rm 2D}\right|$ “radial”
direction as a function of distance from focus $d$. This PSD has a fringe
pattern; the wavevector of the first minimum exceeds the maximum imaged
wavevector only near the image’s focus at $d=0$ ${\mu m}$.
The physical origin of this structure can be understood by turning to the
paraxial wave equations [Eqs. (7) and (8)], and by first studying a single
absorber at ${\bf r}=0$ illuminated by a plane wave $E_{0}^{\prime}({\bf
r}_{\rm 2D},0^{-})=E_{0}$. Equation (8) shows that a thin absorber simply
changes the amplitude of the field, leaving its phase untouched, and for
simplicity, we assume this absorber has a Gaussian profile in the ${\bf
e}_{x}$-${\bf e}_{y}$ plane with width $w_{0}$. Thus the electric field just
following the absorber is changed by $\delta E^{\prime}({\bf r}_{\rm
2D},0^{+})=-\delta E\exp\left[-r^{2}/w_{0}^{2}\right]$, with
$r^{2}=x^{2}+y^{2}$. The propagation of such a gaussian mode by a distance $d$
along ${\bf e}_{z}$ can be solved exactly in the paraxial approximation, and
in the spectral basis this is
$\displaystyle\delta\tilde{E}^{\prime}({\bf k}_{\rm 2D},d)$
$\displaystyle=-\pi w^{2}_{0}\delta E\exp\left[-\frac{w_{0}^{2}k_{\rm
2D}^{2}}{4}\left(1+\frac{2i}{w_{0}^{2}k_{0}}d\right)\right].$
The total field from an absorber located at a different location ${\bf r}_{0}$
in the ${\bf e}_{x}$-${\bf e}_{y}$ plane simply acquires an overall phase
factor $\exp\left[-i{\bf k}_{\rm 2D}\cdot{\bf r}_{0}\right]$. We now compute
the experimentally relevant optical depth by taking the reverse Fourier
transform of the full electric field, computing the intensity, then the
optical depth, and taking the Fourier transform to obtain (retaining terms of
order $\delta E/E_{0}$)
$\displaystyle\widetilde{\rm OD}$ $\displaystyle=\frac{2\pi w_{0}^{2}\delta
E}{E_{0}}\exp\left(-\frac{w_{0}^{2}k_{\rm
2D}^{2}}{4}\right)\cos\left(\frac{k_{\rm 2D}^{2}d}{2k_{0}}\right),$ (13)
with the same overall phase factor depending on the initial position.
Averaging over $N$ randomly placed absorbers therefore gives an overall signal
scaling as $\sqrt{N}$ with a random overall phase. Taking the magnitude
squared gives the PSD
$\displaystyle{\rm PSD}_{\rm thin}$ $\displaystyle={\rm
PSD}_{0}\times\cos^{2}\left(\frac{k_{\rm 2D}^{2}d}{2k_{0}}\right),$
with
$\displaystyle{\rm PSD}_{0}$ $\displaystyle=N\left(\frac{2\pi w_{0}^{2}\delta
E}{E_{0}}\right)^{2}\exp\left(-\frac{w_{0}^{2}k_{\rm 2D}^{2}}{2}\right).$
This quantity has zeros located at $k_{\rm zero}[n]=\sqrt{2\pi(n+1/2)k_{0}/d}$
for integer $n$. In our numerical simulation, the minima follow the functional
form $k_{\rm zero}[n]=A[n]|d|^{-1/2}$ as shown by the dotted lines in Fig. 2,
with $A[0]\approx 5.08$ and $A[1]\approx 8.76$ for the first and second zeros:
the expected values for $A[n]$. Thus, for this thin apparently structureless
system, fringes in the PSD allow us to identify the focal plane.
Figure 2: (a) Spatial PSD of the intensity produced by $1\,\mathrm{\mu m}$
thick layer of randomly distributed scatterers showing that fringes diverge in
focus. (b) PSD produced by a $100\,\mathrm{\mu m}$ thick sheet of random
columnar scatterers. (c) PSD produced by a $100\,\mathrm{\mu m}$ thick sheet
of random scatterers. The dotted lines are functional forms of the lowest
curved-fringes, and in each case $d$ is measured from the objects’ center.
To demonstrate the technique of finding optimal focus of an extended object,
we now consider a second disordered scattering potential with a columnar
structure, now $100\,\mathrm{\mu m}$ thick, i.e., $\chi({\bf r})=ig(x,y)$ for
$z\in(-50\,\mathrm{\mu m},50\,\mathrm{\mu m})$, where again $g(x,y)\geq 0$ is
a Poisson distributed random variable. This object’s PSD is plotted as a
function of distance $d$ from its center in Fig. 2b; in addition to the same
fringe pattern as for the $1\,\mathrm{\mu m}$ thick case, the PSD now vanishes
at specific spatial frequencies independent of $d$. To model this, we note
that the absorbers can now be located at a distance $z$ from the symmetry
plane, so in Eq. (13), we replace $d\rightarrow d-z$ and integrate $z$ from
$-\delta z/2$ to $\delta z/2$, which ultimately gives the PSD
$\displaystyle{\rm PSD}_{\rm col}$ $\displaystyle={\rm
PSD}_{0}\times\cos^{2}\left(\frac{k_{\rm 2D}^{2}d}{2k_{0}}\right){\rm
sinc}^{2}\left(\frac{k_{\rm 2D}^{2}\delta z}{4k_{0}}\right).$
This predicts the appearance of additional zeros located at $k^{\prime}_{\rm
zero}[m]=\sqrt{4\pi mk_{0}/\delta z}$ for non-zero integer $m$ (this is an
artifact of the box-like density distribution of atoms, and would be greatly
softened in real systems where the density drops smoothly to zero). In our
example, the lowest order horizontal fringes is located at $k^{\prime}_{\rm
zero}[1]=1.00\,\mathrm{\mu m^{-1}}$. Here again, we easily determine the
optimal focus, $d=0\ {\mu m}$, from the diverging curved-fringes.
Next, we consider a scattering potential fully disordered in 3D, again with a
$100\,\mathrm{\mu m}$ thickness, i.e., $\chi({\bf r})=ig(x,y,z)$ for
$z\in(-50\,\mathrm{\mu m},50\,\mathrm{\mu m})$, where $g(x,y,z)\geq 0$ is a
Poisson distributed random variable. In this case, the independent random
scatterers along imaging direction causes the PSD to rapidly loose structure
with increasing $k_{\rm 2D}$ (see Fig. 2c). Here too, our random scatter model
can be applied, giving
$\displaystyle{\rm PSD}_{\rm rnd}$ $\displaystyle={\rm
PSD}_{0}\times\left[\frac{\cos\left(k_{\rm 2D}^{2}d/k_{0}\right){\rm
sinc}\left(k_{\rm 2D}^{2}\delta z/2k_{0}\right)+1}{2}\right].$
This reduces to our earlier result when $\delta z\rightarrow 0$ for a thin
system and shows that, while the same fringes exist, they are rapidly
attenuated for larger spatial frequencies, where the signal approaches a
constant background value. However, in principle the curved-fringes still
allow the optimal focus to be identified.
## III Optimal focusing of Elongated Bose-Einstein Condensates
Using on our model, we now consider absorption imaged BECs and implement the
technique presented in previous section to find the optimal focus.
Figure 3: (a) Absorption imaged elongated BEC with density fluctuations. (b)
1D PSD of column density along weakly trap direction ${\bf e}_{x}$ as a
function of $t_{\rm TOF}$. (c) Values of $k$ where the 1D PSD is minimum. The
two lowest such $k$-fringes are depicted. Symbols denote the fringe locations
extracted from (b) plotted along with Lorentzian fits (dotted lines),
determining the optimal focus. The solid curves depict theoretical functional
forms for the two lowest order fringes. In (b) and (c) the dashed line marks
$k=k^{\prime}_{\rm zero}[1]=0.82\,\mathrm{\mu m^{-1}}$ for our condensate
thickness of $150\,\mathrm{\mu m}$; the dotted line marks $k=k^{\prime}_{\rm
zero}[1]/\sqrt{2}$, below which the ACF of the focused images reliably
reflects the ACF $n_{\rm 2D}({\bf r}_{\rm 2D})$.
We prepared $N=7\times 10^{5}$ atom 87Rb Bose-Einstein condensates in the
$\left|5\mbox{S}_{1/2},F=1,m_{F}=0\right\rangle$ electronic ground state in a
crossed-dipole trap with frequencies
$\omega_{x,y,z}=2\pi\times(3.1,135,135)\,\mathrm{Hz}$. In situ, the BECs were
javelin shaped owing to the extremely anisotropic confining potential. After a
17 ms to 21 ms TOF, we repumped into the $f=2$ manifold, and resonantly imaged
on the $\left|5\mbox{S}_{1/2},f=2,m_{F}=2\right\rangle$ to
$\left|5\mbox{P}_{3/2},f=3,m_{F}=3\right\rangle$ cycling transition with a
$\lambda\approx 780.2\,\mathrm{nm}$ probe laser.
The imaging system consisted of a CCD camera and two pairs of lenses
functioning as a compound microscope, magnifying the intensity pattern at the
object by a factor of $\approx 6$ at the image plane. The first pair of
objective lenses, with effective focal length (efl)
$\hat{f_{1}}=53.6\,\mathrm{mm}$, collimated the light diffracted by the cloud
and were separated by a distance $D=\hat{f_{1}}+\hat{f_{2}}$ from a second
pair of lenses with a $\hat{f_{2}}=325\,\mathrm{mm}$ efl. The resulting $0.23$
numerical aperture implies that a $10.6\,\mathrm{\mu m}$ diffraction-limited
spot on our CCD sensor is larger than its $5.6\,\mathrm{\mu m}$ pixel size.
The associated $1.7\,\mathrm{\mu m}$ spot-size on the cloud gives a $d_{\rm
dof}=18.6\,\mathrm{\mu m}$ depth of field in our imaging system Inoue and
Oldenbourg (1995).
Instead of varying the distance from focus by physically moving imaging lenses
or the CCD, we changed the time during which the BEC fell along ${\bf e}_{z}$
and obtained absorption images with TOF times $t_{\rm TOF}$ from
$17.0\,\mathrm{ms}$ to $21.0\,\mathrm{ms}$. At these TOFs, the condensates’
radii were $R_{y,z}\approx 75(5)\,\mathrm{\mu m}$ and $R_{x}\approx
210(10)\,\mathrm{\mu m}$. Initially, the cloud was elongated in the harmonic
trap with aspect ratio $43$ to 1. The initial $43:1$ aspect ratio was reduced
to $2.65:1$ after TOF, and the transverse size of the cloud exceeded the
imaging depth of field by a factor of $8$.
Figure 3b shows the 1D PSD of the atoms’ corrected optical depths along ${\bf
e}_{z}$, which is directly related to the absorption intensity through Eq.
(11). The fluctuations in the BEC’s density distribution behave like the
randomly modulated $\chi({\bf r})$ in our example systems, creating a
recurring fringe pattern in the PSD spectrum as obtained in Fig. 3c. The
fringes are quite pronounced for quasi one-dimensional BECs, where initial
phase fluctuations map into pancake-shaped density fluctuations arrayed along
the initially long axis after TOF Dettmer et al. (2001). Despite the decreased
contrast at high spatial frequencies due to the BEC extent along ${\bf
e}_{z}$, we clearly observe fringes curving as a function of $t_{\rm TOF}$ in
Fig. 3c. This allows us to determine the optimal focus of the system.
From the above experimental data, we fit the two lowest order fringes to
$k_{m}\left[(d-z_{0})^{2}/\delta z^{2}+1\right]^{-1/4}$, a peaked function
with the expected $d^{-1/2}$ behavior away from $z_{0}$. The fits give an
optimal focus location of $z_{0}[0]=1836(2)\,\mathrm{\mu m}$ using the zeroth
order fringe or of $z_{0}[1]=1837(2)\,\mathrm{\mu m}$ using the first order
fringe. These values correspond to a TOF of $19.36(1)\,\mathrm{ms}$. We are
thus able to determine the optimal focus within $\approx 2\,\mathrm{\mu m}$ or
equivalently $\approx 10\,\mathrm{\mu s}$ in TOF. Comparing the experimental
data to the theoretical forms, we notice that the fringes are slightly
asymmetrical with their locations slightly below theoretical ones for larger
TOF. Based on our simulations, this likely results from the $z$ dependent
magnification of our imaging system, which changes by about 10% as the atoms
fall from 1420 $\mu{\rm m}$ to 2150 $\mu{\rm m}$ (17 ms to 21 ms TOF).
## IV Summary
We presented a systematic method to bring clouds of ultracold atoms,
particularly initially elongated BECs, into an optimal focus. The density
fluctuations in the BECs after TOF acted like random scatterers, creating
diffraction pattern which changed predictably as a function of distance from
the optimal focus. Using TOF absorption imaging, we demonstrated this method,
pinpointing the optimal focus of the BEC to within 2 $\mu$m for a
$150\,\mathrm{\mu m}$ thick BEC. This robust technique is easily implemented,
requires no hardware changes, and uses a minimum of computation.
###### Acknowledgements.
We thank F. E. Becerra, A. Hu, and W. D. Phillips for a careful reading of the
manuscript. We acknowledge the financial support from the NSF through the
Physics Frontier Center at JQI, and the ARO with funds from both the
Atomtronics MURI and DARPA’s OLE Program.
## References
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|
arxiv-papers
| 2013-09-19T14:37:27 |
2024-09-04T02:49:51.169761
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Andika Putra and Daniel L. Campbell and Ryan M. Price and Subhadeep De\n and I. B. Spielman",
"submitter": "Andika Putra",
"url": "https://arxiv.org/abs/1309.5000"
}
|
1309.5049
|
∎
11institutetext: R. Liao 22institutetext: B. Bellalta 33institutetext: T.C.
Minh 44institutetext: J. Barcelo 55institutetext: M. Oliver 66institutetext:
Department of Information and Communication Technologies
Universitat Pompeu Fabra, Barcelona, 08018, Spain
66email: {ruizhi.liao, boris.bellalta, trang.cao, jaume.barcelo,
miquel.oliver}@upf.edu
# Uni-MUMAC: A Unified Down/Up-link MU-MIMO MAC Protocol for IEEE 802.11ac
WLANs
Ruizhi Liao Boris Bellalta Trang Cao Minh
Jaume Barcelo Miquel Oliver
###### Abstract
Due to the dominance of the downlink traffic in Wireless Local Area Networks
(WLANs), a large number of previous research efforts have been put to enhance
the transmission from the Access Point (AP) to stations (STAs). The downlink
Multi-User Multiple-Input Multiple-Output (MU-MIMO) technique, supported by
the latest IEEE amendment-802.11ac, is considered as one of the key
enhancements leading WLANs to the Gigabit era. However, as cloud uploading
services, Peer-to-Peer (P$2$P) and telepresence applications get popular, the
need for a higher uplink capacity becomes inevitable.
In this paper, a unified down/up-link Medium Access Control (MAC) protocol
called Uni-MUMAC is proposed to enhance the performance of IEEE 802.11ac WLANs
by exploring the multi-user spatial multiplexing technique. Specifically, in
the downlink, we implement an IEEE 802.11ac-compliant MU-MIMO transmission
scheme to allow the AP to simultaneously send frames to a group of STAs. In
the uplink, we extend the traditional one round channel access contention to
two rounds, which coordinate multiple STAs to transmit frames to the AP
simultaneously. $2$-nd round Contention Window ($CW_{\text{2nd}}$), a
parameter that makes the length of the $2$-nd contention round elastic
according to the traffic condition, is introduced. Uni-MUMAC is evaluated
through simulations in saturated and non-saturated conditions when both
downlink and uplink traffic are present in the system. We also propose an
analytic saturation model to validate the simulation results. By properly
setting $CW_{\text{2nd}}$ and other parameters, Uni-MUMAC is compared to a
prominent multi-user transmission scheme in the literature. The results
exhibit that Uni-MUMAC not only performs well in the downlink-dominant
scenario, but it is also able to balance both the downlink and uplink
throughput in the emerging uplink bandwidth-hungry scenario.
###### Keywords:
MAC MU-MIMO down/up-link IEEE 802.11ac WLANs
## 1 Introduction
IEEE 802.11 Wireless Local Area Networks (WLANs) is becoming an indispensable
part of our life, at homes and working places. Due to the problems, such as
frame collisions and protocol overheads, the throughput of WLANs is
significantly lower than the raw data rate of what the Physical (PHY) layer
can achieve ciscoAC . The evolution of Internet traffic is going to exacerbate
this low-throughput problem. The Internet traffic shifts from web browsings
and file transfers to a wide variety of applications, many of which integrate
content-rich files provided by users DBLP:conf/icumt/KihlOLA10 ;
DBLP:journals/telsys/WamserPSHT11 . This shift, mainly driven by the
bandwidth-hungry multimedia applications (e.g., web HDTV, video sharing and
wireless display), demands a performance increase in both downlink and uplink
of WLANs ciscoWhiteP .
Spatial multiplexing is one of the current trends (the spatial diversity and
the frame aggregation are among others) aiming at improving the performance of
wireless systems. IEEE 802.11n 5307322 supports spatial multiplexing in the
point-to-point communication mode (i.e., Single-user MIMO or SU-MIMO). The
point-to-multipoint communication mode, for example, the transmission from the
Access Point (AP) to multiple stations (STAs) (i.e., downlink Multi-user MIMO
or MU-MIMO), is supported by the latest IEEE amendment-802.11ac 6359961 .
However, the uplink MU-MIMO enhancement, which is crucial to mitigate
collisions and to satisfy the performance requirements in the uploading-
intensive scenario, has not been considered by any IEEE standard.
In this paper, we propose a unified down/up-link MU-MIMO Medium Access Control
(MAC) protocol called Uni-MUMAC, which coordinates distributed STAs to exploit
the spatial multiplexing gain to improve the performance of IEEE 802.11ac
WLANs. The main contributions are summarized as follows. 1) Two separate MU-
MIMO MAC protocols, one for the downlink transmission 5733223 and the other
one for the uplink transmission 6314214 , are integrated into a unified MU-
MIMO MAC protocol. Compared to 5733223 and 6314214 , where only one-way
traffic is considered (i.e., the downlink or the uplink), the presence of both
downlink and uplink transmissions has been taken into account. 2) A special
focus is placed at finding the most suitable value of the $2$-nd round
Contention Window ($CW_{\text{2nd}}$) to obtain the highest system throughput,
and the impact of the optimized uplink transmission on the downlink is
discussed. With the optimized $CW_{\text{2nd}}$ and other properly configured
parameters (e.g., the number of aggregated frames and the queue length of the
AP), Uni-MUMAC is then extensively evaluated through simulations in the
downlink-dominant and the down/up-link balanced traffic scenarios in IEEE
802.11ac based WLANs. 3) An analytic model is developed to validate the
simulation results, and a prominent proposal in the literature is implemented
to compare with our scheme.
The rest of the paper is organized as follows. First, Section 2 explores some
of the key MU-MIMO MAC proposals in the literature. Then, Section 3 introduces
the modified frame structure and detailed Uni-MUMAC operating procedures.
After that, Section 4 gives the considered scenarios to evaluate Uni-MUMAC,
the saturation throughput model, simulation results and observations. Finally,
Section 5 concludes the paper and discusses the future research challenges.
## 2 Related Work
Most previous work has put efforts on adjusting MAC parameters or extending
MAC functions to improve the performance of WLANs. In the downlink, the
spatial multiplexing technique has recently gained much attention. To support
it, many proposals in the literature adopt the following MAC procedure. The AP
firstly sends out a modified Request to Send (RTS) containing a group of
targeted STAs, then those listed STAs estimate the channel, add the estimated
Channel State Information (CSI) into the extended Clear to Send (CTS) and send
it back. As soon as the AP receives all successful CTSs, it precodes the
outgoing signals and sends multiple data frames simultaneously.
Cai et al. in conf/globecom/CaiSZSMW08 propose a distributed MU-MIMO MAC
protocol that modifies RTS and CTS frames to estimate the channel, based on
which, the AP is able to concurrently transmit frames to multiple STAs.
Kartsakli et al. in DBLP:conf/icc/KartsakliZAV09 consider an infrastructured
WLAN and propose four multi-user scheduling schemes to simultaneously transmit
frames to STAs. The results show that the proposal achieves notable gains
compared to that of the single user case. Gong et al. in
DBLP:conf/globecom/GongPSWM10 propose a modified Carrier Sense Multiple
Access with Collision Avoidance (CSMA/CA) protocol with three different ACK-
replying mechanisms. The authors claim that the proposed protocol can provide
a considerable performance improvement against the beamforming based approach
when Signal-to-noise Ratio (SNR) is high. Zhu et al. in
DBLP:conf/ccnc/ZhuBKAN12 investigate the required MAC modifications to
support downlink MU-MIMO transmissions by focusing on the fairness issue. The
proposed Transmit Opportunity (TXOP) sharing scheme not only obtains a higher
throughput but is also more fair than the conventional mechanism. Cha et al.
in 6214021 compare the performance of a downlink MU-MIMO scheme with a Space
Time Block Coding (STBC) based frame aggregation scheme. The results show that
the former produces a higher throughput than the latter if transmitted frames
are of similar length.
The uplink enhancement is getting more attention as the popularity of P$2$P
and cloud applications increases. In general, there are two broad categories
of uplink MU-MIMO MAC enhancements, namely, the un-coordinated access and the
coordinated access. The former utilizes the MAC random mechanism to decide
which STAs are allowed for data transmissions, while the latter employs the AP
to schedule STAs’ uplink access.
Some of the un-coordinated uplink access schemes are sampled as follows. In
DBLP:conf/wcnc/JinJHS08 , Jin et al. evaluate the performance of uplink MU-
MIMO transmissions in the IEEE 802.11 basic access mode, where the
simultaneous uplink transmissions are on the random access basis and the
channel coefficients of each STA are assumed to be known by the AP. In 4025043
, Zheng et al. present a Distributed Coordination Function (DCF) enhancement
called Two-Round RTS Contention (TRRC) to take advantage of the spatial
domain. The proposed scheme allows STAs to contend for the channel after a
successful RTS is detected. In kutse , Tan et al. present a distributed MAC
scheme called Carrier Counting Multiple Access (CCMA), where a beacon that
contains the uplink access threshold is announced by the AP periodically.
Based on the threshold, STAs count the number of ongoing transmissions by
monitoring preambles, and then decide to contend for the channel or stay idle.
In BabichC10 , Babich et al. investigate the theoretical model of asynchronous
frame transmissions, where a STA is allowed to transmit even if other STAs are
already transmitting.
Some of the coordinated uplink access schemes are overviewed as follows. In
ttanc , Tandai et al. propose a synchronized uplink transmission scheme
coordinated by the AP. On receiving requests from STAs, the AP broadcasts a
pilot-Requesting CTS (pR-CTS) to schedule STAs’ pilot transmissions for
estimating the channel. After obtaining the CSI, the AP sends a Notifying-CTS
(N-CTS) to inform the selected STAs to transmit frames in parallel. In SZhou ,
Zhou et al. propose a two-round channel contention mechanism, which divides
the MAC procedure into two parts, namely, the random access and the data
transmission. The random access terminates when the AP receives a predefined
number of successful RTSs, and then the data transmission follows. In
journals/twc/Zhang10 , Zhang et al. further extends the two contention rounds
to multiple rounds, which enable more STAs to be involved in parallel uplink
transmissions. The proposed protocol can fall-back to the single-round mode
automatically on condition that the traffic is low and the single-round scheme
can provide higher throughput. In 6302115 , Jung et al. present an
asynchronous uplink Multi-Packet Reception (MPR) scheme, where an additional
feedback channel is assumed to be employed by the AP to acknowledge the
successful frame receptions along with other ongoing transmissions.
Only a few work has combined the downlink and the uplink transmissions
together. In DBLP:conf/macom/ShenLSDWZ12 , Shen et al. propose a High
Throughput MIMO (HT-MIMO) MAC protocol, which utilizes frequency signatures to
differentiate simultaneously-received control messages. The proposal works in
the Point Coordination Function (PCF) mode, hence both downlink and uplink
transmissions can be only initiated by the AP. In 4917548 , Jin et al. focus
on the unbalanced throughput problem between downlink and uplink, where a
Contention Window (CW) adjustment scheme and a random piggyback scheme are
proposed to increase the downlink throughput ratio. In DBLP:conf/wcnc/LiAL10 ,
Li et al. propose a multi-user transmission MAC scheme, which supports the
Multi-Packet Transmission (MPT) in the downlink and multiple control frame
receptions (e.g., CTSs or ACKs) in the uplink, while simultaneous data
transmissions from multiple STAs are not considered. Due to the simplicity,
the MAC scheme of DBLP:conf/wcnc/LiAL10 is implemented to compare with our
proposal.
## 3 Uni-MUMAC Operations
Uni-MUMAC is based on the IEEE 802.11 Enhanced Distributed Channel Access
(EDCA), which relies on the CSMA/CA mechanism to share the wireless channel.
EDCA can operate in either the basic access mode or the optional RTS/CTS
handshaking one. In this paper, Uni-MUMAC adopts and extends the RTS/CTS
scheme for the following reasons: 1) The AP can notify the uplink contending
STAs about the number of available antennas by a modified control frame; 2)
The AP can estimate the CSI from the RTS/CTS exchanging process; 3) The
distributed STAs can be synchronized from the exchanging process to transmit
to the AP in parallel.
### 3.1 Frame Structure
#### 3.1.1 PHY Frame Structure
The PHY frame structure of IEEE 802.11ac is shown in Figure 1, where VHT PLCP,
PPDU and MPDU stand for Very High Throughput Physical Layer Convergence
Protocol, PLCP Protocol Data Unit and MAC Protocol Data Unit, respectively. As
shown from the frame structure, PPDU consists of the PHY preamble and MPDUs.
IEEE 802.11ac specifies that all MPDUs must be transmitted in the format of
Aggregated-MPDU (A-MPDU), where aggregated MPDUs are separated by MPDU
delimiters. Before being delivered to the PHY layer, a service field and a
tail field are appended to the A-MPDU. The PHY preamble is formed by $3$
legacy fields for the backward compatibility (i.e., L-STF, L-LTF and L-SIG)
and some newly introduced VHT fields 6359961 6140087 .
Figure 1: PHY frame format of IEEE 802.11ac
IEEE 802.11ac introduces these VHT fields to assist WLANs in obtaining the
high performance. A Group Identifier (Group-ID) field is added in VHT Signal
Field-A (VHT-SIG-A), which is used to inform the targeted STAs about the
followed MU-MIMO transmission, the order and the position of each STA’s
corresponding stream. A complete Group-ID table is created and disseminated by
the AP, and will be recomputed as STAs associate or de-associate to the AP.
Since the number of STAs’ combinations can exceed the available number of
Group-ID in a large basic service set, and the down/up-link channel may be
different, thus, we assume a single Group-ID can reference to multiple
transmission sets along with other PHY preamble features that could be used to
resolve the intended STAs aboul2013managing . In other words, there will be
always at least one proper Group-ID entry that can be mapped to the intended
transmission set.
VHT Long Training Field (VHT-LTF) can contain an orthogonal training sequence
that is known by both the transmitter and the receiver to estimate the MIMO
channel. The number of VHT-LTF fields should not be less than the number of
transmitted spatial streams to precisely estimate the channel. The legacy and
VHT-SIG-A fields adopt the low rate modulation scheme to make the preamble
understandable to all STAs, while the rest VHT fields and A-MPDU are
transmitted using the VHT modulation scheme. In this paper, a single
modulation and coding scheme (MCS), i.e., $16$-QAM with $1/2$, is utilized for
all frames to simplify the simulation, although the extension to various MCS
for different frames and STAs is straightforward. Here, we only introduce the
PHY features that are closely related to the proposed protocol. The readers
please refer to 6359961 for details of other PHY features.
#### 3.1.2 MAC Frame Structure
The control frames of Uni-MUMAC are shown in Figures 2 and 3. In the downlink,
the control frames are MU-RTS, MU-CTS and MU-ACK. MU-RTS keeps the standard
RTS frame structure, because the AP can utilize the Group-ID field of the PHY
frame to notify targeted receivers. MU-CTS and MU-ACK add a transmitter
address field to the original CTS and ACK frames in order to facilitate the AP
to differentiate multiple responding STAs. Note that MU-CTS and MU-ACK
coincidentally have the same frame structure as the standard RTS frame after
adding a transmitter address field to the original CTS and ACK frames.
Figure 2: Frame structure of standard RTS
In the uplink, all frame modifications are limited to the AP side to reduce
STAs’ computing consumption. These modified frames are Ant-CTS (CTS with
antenna information), G-CTS (Group CTS) and G-ACK (Group ACK), as shown in
Figure 3. An antenna information field is added to Ant-CTS, which is broadcast
by the AP to announce the number of available antennas (after one antenna is
occupied in the first contention round) and the start of the $2$-nd contention
round. G-CTS and G-ACK have the identical frame structure, where the receiver
address field is removed and replaced by the Group-ID field in the IEEE
802.11ac PHY frame, while a transmitter address field is added to indicate the
AP address. The G-CTS frame is used to inform STAs the start of the data
transmission, and G-ACK is used to indicate the successful reception of data
frames.
(a) Ant-CTS
(b) G-CTS & G-ACK
Figure 3: Modified frames for uplink transmissions
### 3.2 Successful Downlink Transmissions
Figure 4 shows a successful Uni-MUMAC downlink transmission. Initially, the
channel is assumed busy (B). After the channel has been idle for an
Arbitration Inter Frame Space (AIFS), a random backoff (BO) drawn from CW
starts to count down and is frozen as soon as the channel is detected as busy.
Suppose the AP first wins the channel contention and sends a MU-RTS. Then, the
STAs who are included in Group-ID reply with MU-CTSs sequentially as the
indicated order. Those STAs who are not included in the MU-RTS will set the
Network Allocation Vector (NAV) to defer their transmissions. After a MU-CTS
is received, the AP will measure the channel through the training sequence
included in the PHY preamble, and then uses the estimated CSI to precode the
simultaneously-transmitted frames. As being precoded, the frames destined to
different STAs will not interfere with each other. Finally, STAs send MU-ACKs
at the same time to acknowledge the successful reception of data frames.
Figure 4: A successful Uni-MUMAC downlink transmission
Note that, the uplink channel is assumed to be the same as the downlink one in
this paper. In other words, the implicit CSI feedback, namely, the AP
estimates the channel using the training sequence included in the MU-CTS, is
adopted. The reason is that the explicit CSI feedback will need more computing
capability at STAs and require an extra field with substantial volume in the
MU-CTS to include the measured CSI, which may not be suitable for STAs in some
capacity or power constraint scenarios.
### 3.3 Successful Uplink Transmissions
In the uplink, a standard RTS is sent to the AP by the STA that won the $1$-st
round channel contention. Instead of replying a CTS, an Ant-CTS is broadcast
by the AP with two functions: 1) to notify the STA about the successful
reception of the RTS, and 2) to inform other STAs that the number of available
antennas and the start of the $2$-nd contention round. The STAs who have
frames to send will compete for the available spatial streams in the $2$-nd
contention round. A new random $BO$ ($BO_{\text{2nd}}$) drawn from
$[0,CW_{\text{2nd}}-1]$ starts to count down, and a RTS will be sent if
$BO_{\text{2nd}}$ of a STA reaches $0$. The number of available antennas of
the AP decreases by one each time an uplink RTS is successfully received. The
$2$-nd contention round finishes as: 1) all available antennas of the AP are
occupied or 2) a predefined duration of the $2$-nd contention round elapses in
case there are not enough contending STAs (the maximum duration of the $2$-nd
contention round is set to $CW_{\text{2nd}}$ slots). As soon as the $2$-nd
contention round finishes, a G-CTS is sent by the AP to indicate the readiness
for receiving multiple frames in parallel. The G-CTS frame includes the STAs
who have successfully sent RTSs during both $1$-st and $2$-nd contention
rounds. When the G-CTS is received by the targeted STAs, they are synchronized
to send data frames to the AP. Finally, the AP acknowledges the received data
frames with G-ACK.
An example of a successful uplink transmission is shown in Figure 5, in which
illustrating case, the AP has $3$ antennas, STA $2$ picks $BO_{\text{2nd}}=0$
and STA $3$ picks $BO_{\text{2nd}}=1$ from $[0,CW_{\text{2nd}}-1]$,
respectively.
Figure 5: A successful Uni-MUMAC uplink transmission
It is important to point out that the RTSs sent by STAs in the $2$-nd
contention round could collide with G-CTS sent by the AP. For example, in the
case that the RTS sent by a STA who claims the AP’s last available antenna is
not heard by some STAs (hidden terminals), which therefore believe that the AP
still has available antennas. Then, after a Short Inter Frame Space (SIFS)
interval, the G-CTS sent by the AP and RTSs sent by the hidden STAs would
collide. To avoid this unexpected scenario, STAs are forced to wait for a
Multi-User SIFS interval in the $2$-nd contention round. MU-SIFS is an
interval longer than SIFS but shorter than AIFS, which not only prioritizes
the AP to send the G-CTS, but also avoids STAs to misunderstand MU-SIFS as an
idle channel.
### 3.4 Frame Collisions
Collisions will occur in both $1$-st and $2$-nd contention rounds if more than
one STA choose the same random backoff value. On sending a RTS, EDCA specifies
that the STA has to set a timer according to Equation (1) to receive the
expected CTS, where $T_{\text{CTS}}$ represents the transmission duration of a
CTS frame. If CTS is not received before the timer expires, the STAs who
previously sent RTSs assume that collisions occurred. These RTS-sending STAs
will compete for the channel access after the expiration of the timer. For the
RTS-receiving STAs, none of RTSs can be decoded correctly. Therefore, after
the collision time, the receiving STAs will wait for an Extended Inter Frame
Space (EIFS, as shown in Equation (2)) interval to compete for the channel
access together with those RTS-sending STAs.
As shown in Figure 6 (Ant-CTS and MU-CTSs with dotted lines mean these frames
would be transmitted if there were no collisions), collisions in the $1$-st
contention round include two cases: 1) collisions among STAs; 2) collisions
between STAs and the AP. Since STAs can not differentiate these two cases, the
collision time has to be set according to the duration of the longer frame,
which is $T_{\text{MU-RTS}}$. In addition, the $\text{CTS}_{\text{timer}}$ and
the EIFS interval also have to be extended according to $\text{MU-
CTS}_{\text{timer}}$ (as shown in Equation (3), where $N$ is the number of
AP’s antennas) and Multi-User EIFS (MU-EIFS, as shown in Equation (4)), to
take the scenario that the AP is involved in collisions into account.
Figure 6: Collisions in the $1$-st contention round
$\displaystyle\text{CTS}_{\text{timer}}=\text{SIFS}+T_{\text{CTS}}$ (1)
$\displaystyle\text{EIFS}=\text{SIFS}+T_{\text{CTS}}+\text{AIFS}$ (2)
$\displaystyle\text{MU-CTS}_{\text{timer}}=N\cdot(\text{SIFS}+T_{\text{MU-
CTS}})$ (3) $\displaystyle\text{MU-EIFS}=N\cdot(\text{SIFS}+T_{\text{MU-
CTS}})+\text{AIFS}$ (4)
If collisions occur in the $2$-nd contention round, the colliding STAs will
not be indicated as the receivers in the Group-ID field of G-CTS. Therefore,
only the STAs that have successfully sent RTSs in both contention rounds are
allowed to transmit frames to the AP at the same time (as illustrated in
Figure 7).
Figure 7: RTS collisions in the $2$-nd contention round
### 3.5 Other Considerations
In IEEE 802.11 EDCA, a STA renews its $BO$ if the channel contention was
successful. For the STAs who did not win the contention, the frozen $BO$ is
used for the next contention round. In this paper, $BO$ of the $1$-st
contention round is renewed after collisions in the $1$-st round or if the STA
is the initiator of the two-round process. Although both STA $1$ and STA $2$
participate in the transmission as shown in Figure 7, STA $1$ is considered to
be the initiator. In other words, STA $1$ will have a new random $BO$ in the
followed $1$-st contention round, while STA $2$ will use the frozen $BO$.
It is more straightforward regarding the $BO_{\text{2nd}}$ renewal policy.
Each STA draws a fresh $BO_{\text{2nd}}$ from $CW_{\text{2nd}}$ as soon as a
new $2$-nd contention round starts.
G-CTS is sent out by the AP when the number of available antennas reaches zero
or the duration of the $2$-nd contention round drains. As soon as the Ant-CTS
is sent, the AP sets the G-CTS timer to account for up to $CW_{\text{2nd}}$
slots (as shown in Equation (5)).
$\displaystyle\text{G-CTS}_{\text{timer}}=CW_{\text{2nd}}\cdot(\text{MU-
SIFS}+T_{\text{RTS}})$ (5)
## 4 Performance Evaluation
Uni-MUMAC is evaluated using an analytic model and simulations. The analytic
model is adapted from Bianchi’s saturation throughput model
bianchi2000performance to support MU-MIMO transmissions in both downlink and
uplink. The simulation is implemented in `C++` using the Component Oriented
Simulation Toolkit (COST) library GilChen and the SENSE simulator
chen2005sense .
A single-hop WLAN implementing Uni-MUMAC is considered as shown in Figure 8.
It consists of one AP and $M$ STAs with an error-free channel. The AP employs
an array of $N$ antennas, while each STA has only one antenna. The data frame
has a fixed length of $L$ bits. The parameters used to evaluate Uni-MUMAC are
listed in Table 1.
Figure 8: Down/Up-link Uni-MUMAC transmissions
### 4.1 Saturation Throughput Analysis
Table 1: System parameters Parameters | Values
---|---
Channel Bandwidth | $40$ MHz
Modulation & Coding Scheme | $16$-QAM with $1/2$
Guard Interval | $0.8$ $\mu$s
Queue Length of STA & AP | $Q_{\text{sta}}=50$, $Q_{\text{ap}}=M^{2}$
Frame Length ($L$) | $8000$ bits
MAC Header ($L_{\text{MAC}}$) | $272$ bits
MPDU Delimiter ($L_{\text{delimiter}}$) | $32$ bits
Service Bits ($L_{\text{service}}$) | $16$ bits
Tail Bits ($L_{\text{tail}}$) | $6$ bits
RTS/MU-RTS/MU-CTS/MU-ACK | $160$ bits
Ant-CTS | $120$ bits
G-CTS/G-ACK | $112$ bits
Idle Slot ($\sigma$) | $9$ $\mu$s
SIFS, MU-SIFS and AIFS | $16$, $20$ and $34$ $\mu$s
$CW$ | $32$
AP Antennas ($N$) | $1$, $2$, $4$
No. of Iteration ($N_{\text{iteration}}$) | $100000$
Let $\tau=\frac{2}{(CW+1)}$ be the transmission probability of a node in a
random slot, where $CW$ is the size of the $1$-st round contention window.
Then, the probability that the channel is idle is:
$p_{\text{i}}=(1-\tau)^{M+1}.$ (6)
The probability that the channel sees a successful transmission slot,
$p_{\text{s}}$, is given by:
$\displaystyle
p_{\text{s}}=\binom{M+1}{1}\tau(1-\tau)^{M}=(M+1)\tau(1-\tau)^{M},$ (7)
which accounts for that a single node (either the AP or a STA) successfully
wins the $1$-st round channel contention.
By deducting $p_{\text{i}}$ and $p_{\text{s}}$, the probability that the
channel observes a collision slot, $p_{\text{c}}$, is obtained:
$\displaystyle p_{\text{c}}=1-p_{\text{i}}-p_{\text{s}}.$ (8)
In the saturated condition, a successful downlink transmission always contains
$N$ (the number of AP antennas) data streams. Therefore, the number of bits of
a successful downlink transmission ($N_{\text{b,down}}$) is:
$\displaystyle N_{\text{b,down}}=\alpha\cdot N\cdot N_{\text{f}}\cdot L\cdot
p_{\text{s}},$ (9)
where $\alpha=\frac{1}{M+1}$ is the probability that a transmission is from
the AP, and $N_{\text{f}}$ is the number of aggregated frames in an A-MPDU.
The calculation of the successfully received number of bits of uplink
($N_{\text{b,up}}$) has to account for successful transmissions of both $1$-st
and $2$-nd contention rounds:
$\displaystyle N_{\text{b,up}}=(1-\alpha)\cdot N_{\text{f}}\cdot L\cdot
p_{\text{s}}\cdot\sum_{\text{x}=1}^{N}p_{\text{x\\_ant}}\cdot\text{x},$ (10)
where $p_{\text{x\\_ant}}$ is the probability that x (x $\in[1,N]$) antennas
of the AP have been used for the uplink transmission. In other words, one
antenna has been obtained by a STA in the $1$-st contention round, and x-1
antennas have been successfully obtained by STAs in the $2$-nd contention
round.
The duration of a successful downlink transmission, $T_{\text{s,down}}$, is:
$\displaystyle T_{\text{s,down}}=\text{AIFS}+T_{\text{MU-
RTS}}+N\cdot(T_{\text{MU-CTS}}+\text{SIFS})+T_{\text{A-MPDU}}+T_{\text{MU-
ACK}}+2\cdot\text{SIFS}.$ (11)
An example to calculate the duration of a MU-RTS frame and a data frame using
the system parameters of Table 1 is given in Equation (12).
$T_{\text{PHY}}(N)=36+N\cdot 4$ $\mu$s are the duration of PHY header (the
number of the VHT-LTF fields is proportional to the number of AP antennas
$N$); $L_{\text{service}}$, $L_{\text{tail}}$ and $L_{\text{delimiter}}$ are
the length of the service field, the tail field and the MPDU delimiter;
$L_{\text{DBPS}}$ and $T_{\text{symbol}}$ are the number of data bits in a
symbol and the symbol duration; $N_{\text{f}}$ is the number of aggregated
frames in an A-MPDU; $L_{\text{MU-RTS}}$ and $L_{\text{MAC}}$ are the length
of MU-RTS and the MAC header respectively. More detailed calculation of the
frame duration can be found in 6287486 .
$\displaystyle\begin{cases}T_{\text{MU-
RTS}}=T_{\text{PHY}}(N)+\Big{\lceil}\frac{L_{\text{service}}+L_{\text{MU-
RTS}}+L_{\text{tail}}}{L_{\text{DBPS}}}\Big{\rceil}T_{\text{symbol}}\\\
T_{\text{A-MPDU}}=T_{\text{PHY}}(N)+\Big{\lceil}\frac{L_{\text{service}}+N_{\text{f}}\cdot(L_{\text{MAC}}+L+L_{\text{delimiter}})+L_{\text{tail}}}{L_{\text{DBPS}}}\Big{\rceil}T_{\text{symbol}}\end{cases}$
(12)
The duration of a successful uplink transmission, $T_{\text{s,up}}$, is:
$\displaystyle T_{\text{s,up}}=\text{AIFS}+T_{\text{RTS}}+T_{\text{Ant-
CTS}}+E_{\text{2nd-
slots}}+T_{\text{G-CTS}}+T_{\text{A-MPDU}}+T_{\text{G-ACK}}+4\cdot\text{SIFS},$
(13)
where $E_{\text{2nd-slots}}$ stands for the average duration of the $2$-nd
contention round.
$\displaystyle E_{\text{2nd-slots}}=(T_{\text{RTS}}+\text{MU-
SIFS})\cdot\sum_{\text{k}=1}^{CW_{\text{2nd}}}p_{\text{k\\_Slot}}\cdot\text{k},$
(14)
where $p_{\text{k\\_Slot}}$ is the probability that there are k (k
$\in[1,CW_{\text{2nd}}]$) slots in the $2$-nd contention round.
As a STA can not differentiate if collisions of the $1$-st round are caused by
the AP or other STAs, the collision time has to be set according to the
duration of the longer frame:
$\displaystyle T_{\text{c}}=\text{AIFS}+T_{\text{MU-RTS}}+N\cdot(T_{\text{MU-
CTS}}+\text{SIFS}).$ (15)
The average duration of a channel slot is:
$\displaystyle T_{\text{average}}=\alpha\cdot p_{\text{s}}\cdot
T_{\text{s,down}}+(1-\alpha)\cdot p_{\text{s}}\cdot
T_{\text{s,up}}+p_{\text{c}}\cdot T_{\text{c}}+p_{\text{i}}\cdot\sigma.$ (16)
Equation (17) gives a simple example to calculate $p_{\text{2\\_ant}}$, in
which case, the AP has $2$ antennas and $CW_{\text{2nd}}=2$:
$\displaystyle
p_{\text{2\\_ant}}=\binom{M-1}{1}\frac{1}{CW_{\text{2nd}}}\left(1-\frac{1}{CW_{\text{2nd}}}\right)^{M-2}+\binom{M-1}{1}\frac{1}{CW_{\text{2nd}}}\left(1-\frac{1}{CW_{\text{2nd}}}\right)^{M-2}\cdot
p_{\text{1\\_fail}}.$ (17)
The first part of Equation (17) stands for that only one STA is successful in
the $1$-st slot. The second part represents that only one STA is successful in
the $2$-nd slot, which is conditioned on that the $1$-st slot fails
($p_{\text{1\\_fail}}$, no STAs or more than one STA chooses the $1$-st slot).
Note that the similar condition is not required for the first part, because
the $2$-nd round contention finishes as soon as a STA wins the $1$-st slot
regardless the choices of other STAs of other slots. As $CW_{\text{2nd}}$
increases, the closed form of $p_{\text{2\\_ant}}$ becomes infeasible due to
various combination of conditions for a STA to succeed in different slots.
Therefore, we utilize a semi-analytic algorithm to calculate
$p_{\text{x\\_ant}}$ and $p_{\text{k\\_Slot}}$, the pseudo code of which is
shown in Algorithm 1.
Input: $N$, $M$, $CW_{\text{2nd}}$, $N_{\text{iteration}}$
Output: $p_{\text{x\\_ant}}$, $p_{\text{k\\_Slot}}$
for _$i\leq N_{\text{iteration}}$_ do
for _$\textbf{STA\\_id}=1:M$_ do
$rand\\_value(STA\\_id)=\lfloor rand()*CW_{\text{2nd}}\rfloor$
$\ignorespaces\triangleright$ $rand\\_value$: random value chosen by each STA
$CW2nd\\_Array(rand\\_value(STA\\_id))\texttt{++}$
$\ignorespaces\triangleright$ $CW2nd\\_Array$: no. of times a
$CW_{\text{2nd}}$ value is chosen
end for
$count=1$
$1st\\_success=\textbf{false}$
for _$\textbf{length}=1:CW_{\text{2nd}}$_ do
if _$CW2nd\\_Array(length)\texttt{==}1$_ then
$num\\_sta\\_success(count)\texttt{++}$ $\ignorespaces\triangleright$
$num\\_sta\\_success$: no. of STAs won the $2$-nd round contention if
_$1st\\_success\texttt{==}\textbf{true}$_ then
$num\\_sta\\_success(count\texttt{-}1)\texttt{--}$
end if
if _$count\texttt{==}1$_ then
$1st\\_success=\textbf{true}$
end if
if _$count\texttt{==}N\texttt{-}1$_ then
break
end if
$count\texttt{++}$
end if
if _$length\texttt{==}CW_{\text{2nd}}$_ then
$slot\\_id=CW_{\text{2nd}}$
end if
end for
$num_{\text{slot\\_success}}(slot\\_id)\texttt{++}$
$\ignorespaces\triangleright$ $num_{\text{slot\\_success}}$: no. of slots used
in the $2$-nd round contention
end for
for _$\text{x}=1:N\texttt{-}1$_ do
$p_{\text{(x+1)\\_ant}}=num\\_sta\\_success(\text{x})/N_{\text{iteration}}$
end for
$p_{\text{1\\_ant}}=1-\sum_{x=2}^{N}p_{\text{x\\_ant}}$
for _$k=1:CW_{\text{2nd}}$_ do
$p_{\text{k\\_Slot}}=num_{\text{slot\\_success}}(\text{k})/N_{\text{iteration}}$
end for
Algorithm 1 Algorithm for $p_{\text{x\\_ant}}$ and $p_{\text{k\\_Slot}}$
Finally, the collision probability of a node,
$\displaystyle P_{\text{collision}}=1-(1-\tau)^{M},$ (18)
and down/up-link throughput are derived:
$\displaystyle\begin{cases}S_{\text{down}}=\frac{N_{\text{b,down}}}{T_{\text{average}}}\\\
S_{\text{up}}=\frac{N_{\text{b,up}}}{T_{\text{average}}}.\end{cases}$ (19)
The transmission probability $\tau$, equations (18) and (19) form a non-linear
system, which can be resolved by an iterative numerical technique kumar2005new
.
### 4.2 System Performance against $CW_{\text{2nd}}$
In this sub-section, the performance of Uni-MUMAC is evaluated by increasing
$CW_{\text{2nd}}$, with the goal to find a suitable $CW_{\text{2nd}}$ value
that maximizes the system performance. Two traffic conditions are considered:
1) the saturated one, as shown in Figure 9, and 2) the non-saturated one, as
shown in Figure 10. The saturated condition means that both the AP and STAs
always have frames to transmit. Obviously, there is no $2$-nd round channel
access when the AP has $1$ antenna, which is why the results keep constant as
$N=1$. Note that the plots include both analysis and simulation results in the
saturated condition, while the plots include only simulation results of the
non-saturated condition.
As shown in Figure 9, when the WLAN is saturated (i.e., both downlink and
uplink are saturated), $CW_{\text{2nd}}$ has very small impact on the downlink
throughput (AP’s throughput). However, for the uplink, the importance of
choosing an appropriate $CW_{\text{2nd}}$ is observed. For example, the uplink
throughput (STAs’ throughput) approaches its maximum when
$CW_{\text{2nd}}\in[8,12]$ as $M=8$ (Figure 9(a)) and when
$CW_{\text{2nd}}\in[12,16]$ as $M=15$ (Figure 9(b)).
(a) $M=8$
(b) $M=15$
Figure 9: Saturated throughput against $CW_{\text{2nd}}$
In the non-saturated condition, we set the traffic load for each STA and the
AP to $1.4$ Mbps and $11.2$ Mbps, respectively. In Figure 10(a), the downlink
throughput ($N=2$ and $4$) obtains the highest value when
$CW_{\text{2nd}}\in[4,8]$, and then decreases as $CW_{\text{2nd}}$ keeps
increasing. The reason is that the continuous increase of $CW_{\text{2nd}}$
leads to longer uplink transmissions that harm the downlink ones. Figure 10(b)
shows that the average delay increases as $CW_{\text{2nd}}$ increases. Note
that, the average delay remains at a relatively low level when the system is
in the non-saturated condition, for example, the average delay of STAs when
$CW_{\text{2nd}}\in[4,34]$ and the average delay of the AP when $N=4$ and
$CW_{\text{2nd}}\in[4,8]$. However, the average delay of the AP ($N=4$)
increases sharply as the downlink traffic approaches saturation.
(a) $M=8$, STA $1.4$ Mbps, AP $11.2$ Mbps
(b) $M=8$, STA $1.4$ Mbps, AP $11.2$ Mbps
Figure 10: Non-saturated throughput & Average delay
It is also observed that the downlink throughput, as the network becomes
saturated, is much lower than the uplink one. The reasons are as follows.
First, the AP bottle-neck effect. It is due to the fact that the AP manages
all traffic to and from STAs in a WLAN, while it has the same probability to
access the channel as the STAs due to the random backoff mechanism of CSMA/CA.
In addition, the inherently high traffic load at the AP results in that the
downlink is saturated in most of the time. Thirdly, a favorable value of
$CW_{\text{2nd}}$ for the uplink does not mean the same benefit to the
downlink. For example, as shown in the Figure 9, the uplink obtains the
highest throughput when $CW_{\text{2nd}}$ is set approximately to $M$
($CW_{\text{2nd}}\approx M$), while the downlink transmission prefers a value
of $CW_{\text{2nd}}$ as small as possible.
In order to mitigate the AP bottle-neck effect and compensate the downlink
disadvantage when STAs choose a big $CW_{\text{2nd}}$, we set the maximum
number of frames that the AP can aggregate in an A-MPDU to $M$
($N_{\text{f}}\leq M$), while keeping the number of frames aggregated by each
STA to $1$ in the following simulations. Also, the queue length of the AP is
set to quadratically increase with the number of STAs ($Q_{\text{ap}}=M^{2}$)
to statistically guarantee that there are enough frames destined to different
STAs 6287486 .
In Figures 11 and 12, the performance of Uni-MUMAC is evaluated in the same
condition as done in Figures 9 and 10 except that the network adopts the new
frame aggregation scheme (AP’s $N_{\text{f}}\leq M$, STA’s $N_{\text{f}}=1$)
and the new queue length ($Q_{\text{ap}}=M^{2}$, $Q_{\text{sta}}=50$). The
results show that Uni-MUMAC manages to avoid the extremely low downlink
throughput when the system is saturated (Figure 11) and keeps the downlink
transmission always in the non-saturation area (Figure 12(a)), which is not
achieved in Figure 10(a). The average delay of the AP (Figure 12(b)) is much
lower compared to that of the AP in Figure 10(b), which is because the system
remains in the non-saturated condition by employing the frame aggregation
scheme.
(a) $M=8$
(b) $M=15$
Figure 11: Saturated throughput when AP aggregates frames
The results from Figure 11 also show that the system can roughly obtain the
maximum performance when $CW_{\text{2nd}}\in[M-4,M+4]$. For example, in the
case that the AP has $4$ antennas, the system throughput (AP+STA) reaches its
maximum when $CW_{\text{2nd}}\in[6,8]$ as $M=8$ and
$CW_{\text{2nd}}\in[12,16]$ as $M=15$, respectively. Therefore, the optimum
value of $CW_{\text{2nd}}$ is fixed to $M$ in the following simulations.
(a) $M=8$, STA $1.4$ Mbps, AP $11.2$ Mbps
(b) $M=8$, STA $1.4$ Mbps, AP $11.2$ Mbps
Figure 12: Non-saturated throughput & Average delay when AP aggregates frames
### 4.3 System Performance against $M$
In this sub-section, the performance of Uni-MUMAC is evaluated against the
number of STAs in the downlink-dominant and the down/up-link balanced traffic
scenarios, where $M$ is increased from $1$ to $35$, the maximum number of
frames aggregated at the AP is set to $M$ and the $2$-nd round Contention
Window is also set to $M$. The two traffic scenarios are specified as follows.
1. 1.
Downlink-dominant: This is the traditional WLAN traffic scenario, where the AP
manages a much heavier traffic load compared to that of STAs. Therefore, the
traffic load of the AP is set to be $4$ times higher than that of each STA.
For instance, if the traffic load of a STA is $0.8$ Mbps and there are $5$
STAs, the traffic load of the AP will be $4\cdot 0.8\cdot 5=16$ Mbps.
2. 2.
Down/up-link balanced: This is one of WLAN traffic types that not only
includes P$2$P applications, which have already been around for some years,
but also includes those emerging content-rich file sharing and video calling
applications. Therefore, the traffic load of the AP is set to be the same as
that of each STA. In this case, if there are $5$ STAs, and each STA has $0.8$
Mbps traffic load, the traffic load of the AP will be $0.8\cdot 5=4$ Mbps.
The multi-user MAC scheme (LI-MAC) proposed by Li et al. in
DBLP:conf/wcnc/LiAL10 is implemented and used as a reference (named as
AP/STA-LI in the legend) to compare with Uni-MUMAC. For fair comparison, LI-
MAC and Uni-MUMAC adopt the same configuration parameters (as shown in Table
1). The key features of LI-MAC and Uni-MUMAC are illustrated in Table 2.
Table 2: Key features of LI-MAC and Uni-MUMAC
MAC Schemes Modification Downlink Uplink LI-MAC DBLP:conf/wcnc/LiAL10 MU-
RTS/CTS Multi-packet + Parallel-control-frame TX One-round single-packet TX
Uni-MUMAC MU-RTS/CTS/ACK/SIFS G-CTS/ACK, Ant-CTS Multi-packet + Sequential-
control-frame TX Two-round Multi-packet TX
Figure 13(a) shows the throughput by increasing the number of STAs in the
downlink-dominant traffic scenario. It is with clear advantage to employ a
higher number of antennas at the AP. The downlink throughput is much higher
than the uplink one before the system gets saturated. The reasons for that are
twofold: 1) the AP traffic load is inherently higher than that of STAs, and 2)
the AP adopts the frame aggregation scheme. As the system becomes saturated,
the throughput of both downlink and uplink decreases as $M$ increases.
As shown in Figure 13(a), the uplink throughput of LI-MAC ($N=4$) is the same
as that of Uni-MUMAC ($N=1$), which is because LI-MAC adopts the baseline DCF
in the uplink. As the uplink throughput approaches saturation ($M=15$), the
downlink throughput of LI-MAC starts to decrease. The downlink throughput of
Uni-MUMAC can achieve higher gains when the network is not saturated, which is
because the proposed $2$-nd round transmission increases the uplink
transmission efficiency, and therefore decreases the number of AP’s channel
contenders. However, as the number of STAs further increases, where both
up/down-link saturate, LI-MAC outperforms Uni-MUMAC, which is because Uni-
MUMAC suffers a high collision rate in the 2-nd round that prolongs the 2-nd
round duration. However, it is important to point out that neither LI-MAC or
Uni-MUMAC is able to work sustainably in the saturated condition.
(a) Downlink-dominant: STA $0.8$ Mbps, AP $3.2$ Mbps
(b) P$2$P Scenario: STA $0.8$ Mbps, AP $0.8$ Mbps
Figure 13: Throughput against $M$
Figure 13(b) shows the throughput against $M$ in the down/up-link balanced
traffic scenario. As expected, Uni-MUMAC achieves the balanced downlink and
uplink throughput. This is because the AP and STAs are set to have the same
traffic load, and more importantly, the frame aggregation scheme (AP’s
$N_{\text{f}}\leq M$, STA’s $N_{\text{f}}=1$) counteracts the STAs’ collective
advantage on the channel access.
Comparing with Uni-MUMAC, the downlink throughput of LI-MAC achieves better
performance when the uplink is saturated, which is because the duration of
collisions in the uplink of LI-MAC is much shorter than that of Uni-MUMAC.
However, the drawback is that LI-MAC has a big throughput gap between the AP
and STAs, which does not satisfy the traffic requirements of the considered
scenario.
Figure 14 shows the average delay against $M$. Both downlink and uplink delays
increase with $M$, and grow significantly as the downlink or the uplink
traffic approaches the saturation. After the system gets saturated, the
average delay becomes steady. It is worth pointing out that the average delay
of STAs is higher than that of the AP when $M$ becomes bigger. The reason for
that is that the transmission duration of the AP gets longer as $M$ increases
(due to the frame aggregation scheme), which makes STAs waiting longer to
access the channel.
(a) Downlink-dominant: STA $0.8$ Mbps, AP $3.2$ Mbps
(b) P$2$P Scenario: STA $0.8$ Mbps, AP $0.8$ Mbps
Figure 14: Average delay against $M$
Figure 15 shows the $1$-st round collision probability increases with $M$ and
converges when the system becomes saturated, which confirms the down/up-link
saturation trend as discussed in Figures 13 and 14. It is interesting to note
that the collision probability of STAs is higher than that of the AP when the
system is non-saturated. The reason for that is a STA transmits less
frequently than the AP in the non-saturated condition, which results in a
lower conditional collision probability for the AP. It can be clearly
explained by Equation 20, where $p_{\text{ap}}$ and $\tau_{\text{ap}}$
($p_{\text{sta}}$ and $\tau_{\text{sta}}$) are the $1$-st round collision
probability and the transmission probability of the AP (or a STA) in the non-
saturated condition:
(a) Downlink-dominant: STA $0.8$ Mbps, AP $3.2$ Mbps
(b) P$2$P Scenario: STA $0.8$ Mbps, AP $0.8$ Mbps
Figure 15: $1$-st round collision probability against $M$
$\displaystyle\begin{cases}p_{\text{ap}}=1-(1-\tau_{\text{sta}})^{M}\\\
p_{\text{sta}}=1-(1-\tau_{\text{sta}})^{M-1}\cdot(1-\tau_{\text{ap}}).\end{cases}$
(20)
Figure 16 shows the $2$-nd round collision probability against $M$. It is
clear that the $2$-nd round collision probability is higher when the system
traffic load is higher. In the low number of STAs area, the $2$-nd round
collision probability when the AP has $2$ antennas is sometimes lower than
that when the AP has $4$ antennas. The reason is that, a higher number of
antennas at the AP usually means a longer duration of the $2$-nd contention
round, which increases the chances of collisions in the $2$-nd round. For
example, in a case that the AP employs $2$ antennas, the $2$-nd contention
round finishes as soon as a STA successfully wins the still-available antenna
of the AP; while in a case that the AP employs more than $2$ antennas, the
$2$-nd contention round continues, therefore increasing the $2$-nd round
collision probability.
Figure 16: $2$-nd round collision probability against $M$
## 5 Conclusions & Future Research Challenges
In this paper, a unified MU-MIMO MAC protocol called Uni-MUMAC, which supports
both MU-MIMO downlink and uplink transmissions for IEEE 802.11ac WLANs, is
proposed. We evaluate it through an analytic model and simulations. A
prominent MAC scheme from the literature is implemented and compared with Uni-
MUMAC.
By analyzing the simulation results, we observe that the $2$-nd round
Contention Window $CW_{\text{2nd}}$, which is tuned to optimize the uplink
transmission, is however not bringing the same benefit to the downlink one. An
adaptive frame aggregation scheme and a queue scheme are applied at the AP to
offset this disadvantage. By properly setting the aforementioned parameters,
the results show that a WLAN implementing Uni-MUMAC is able to avoid the AP
bottle-neck problem and performs very well in both the traditional downlink-
dominant and emerging down/up-link balanced traffic scenarios. The results
also show that a higher system capacity can be achieved by employing more
antennas at the AP.
Uni-MUMAC gives us insight about the interaction of down/up-link transmissions
and how different parameters that control the system can be tuned to achieve
the maximum performance. Based on the study of this paper, we considered the
following aspects as the future research challenges or next steps for Uni-
MUMAC.
1. 1.
Adaptive Scheduling Scheme: As discussed in the paper, a parameter that
optimizes the uplink could be unfavorable to the downlink. Therefore, an
adaptive scheduling algorithm that takes several key parameters into account
and compensates those STAs whose interests are harmed would play a significant
role on obtaining the maximum performance while maintaining the fairness. As
implied from the results, these parameters include: the size of A-MPDU, the
queue length, the spatial-stream/frame allocation, the number of
nodes/antennas, and other key parameters that control down/up-link
transmissions.
2. 2.
Traffic Differentiation: Another future research challenge is to provide new
traffic differentiation capability in the uplink in addition to the one
defined in IEEE 802.11e amendment 11e . The new traffic differentiation should
be able to limit the number of STAs that can participate in the $2$-nd
contention round to reduce $2$-nd round collisions. A possible solution could
be to create a table at the AP with information about the priority of each
traffic flow and the queue length of each STA, and then to utilize this table
to control the $2$-nd contention round.
3. 3.
Multi-hop Mesh Networks: There are more challenges that need to be considered
in designing MAC to operate in multi-hop wireless networks. First, the hidden-
node problem. It is still an open challenge to find mechanisms that
efficiently solve the collisions caused by hidden nodes. A collision-free
scheme proposed in barcelo2012distributed or the handshake based coordinated
access could be a starting point to combat the hidden-node collisions in
wireless mesh networks. Secondly, due to the heterogeneity of mesh nodes
(e.g., different number of antennas at nodes), MAC protocols for wireless mesh
networks need to be designed with the capability of swiftly switching among
MU-MIMO, SU-MIMO, multi-packet and single-packet transmission schemes.
Thirdly, MAC and routing protocols need to be jointly designed. There could be
multiple destinations involved in a MU-MIMO transmission, and some
destinations could be out of the one-hop transmitting range, in which case,
routing strategies should be able to forward multiple packets to different
nodes in parallel.
###### Acknowledgements.
This work has been supported by the Spanish Government and the Catalan
Government under projects TEC2012-32354 (Plan Nacional I+D), CSD2008-00010
(Consolider-Ingenio Program) and SGR2009#00617.
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|
arxiv-papers
| 2013-09-19T16:49:16 |
2024-09-04T02:49:51.179424
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Ruizhi Liao, Boris Bellalta, Trang Cao Minh, Jaume Barcelo, Miquel\n Oliver",
"submitter": "Ruizhi Liao",
"url": "https://arxiv.org/abs/1309.5049"
}
|
1309.5064
|
11institutetext: US Naval Research Laboratory-Code 6792, Nonlinear System
Dynamics Section, Plasma Physics Division, Washington, DC 20375
Bloomberg School of Public Health-Johns Hopkins University, 615 N. Wolfe
Street, Baltimore, MD 21205
# Capturing pattern bi-stability dynamics in delay-coupled swarms
L. Mier-y-Teran-Romero and I. B. Schwartz 1122111122
###### Abstract
Swarms of large numbers of agents appear in many biological and engineering
fields. Dynamic bi-stability of co-existing spatio-temporal patterns has been
observed in many models of large population swarms. However, many reduced
models for analysis, such as mean-field (MF), do not capture the bifurcation
structure of bi-stable behavior. Here, we develop a new model for the dynamics
of a large population swarm with delayed coupling. The additional physics
predicts how individual particle dynamics affects the motion of the entire
swarm. Specifically, (_1_) we correct the center of mass propulsion physics
accounting for the particles’ velocity distribution; (_2_) we show that the
model we develop is able to capture the pattern bi-stability displayed by the
full swarm model.
###### pacs:
05.45.-a
###### pacs:
89.75.Kd
###### pacs:
87.23.Cc
Nonlinear dynamics Pattern formation in complex systems Population dynamics
## 1 Introduction
Recently, much attention has been given to the study of interacting multi-
agent, particle or swarming systems in various natural [1, 2, 3, 4] and
engineering [5, 6] fields. These multi-agent swarms can self-organize and form
complex spatio-temporal patterns even when the coupling between agents is
weak. Many of these investigations have been motivated by a multitude of
biological systems such as schooling fish, swarming locusts, flocking birds,
bacterial colonies, ant movement, etc. [7, 8, 9, 10, 11, 12, 13, 14], and have
also been applied to the design of systems of autonomous, communicating robots
or agents [15, 16, 17] and mobile sensor networks [18, 19]. The excellent
overviews [20, 21] discuss the diverse biological contexts where swarming
occurs, as well as different modeling approaches. We note that in spite of all
these investigations, understanding how swarming patterns self-organize, as
well as predicting their stability are still very much open problems.
A number of different mathematical modeling techniques have been applied to
investigate aggregating particle systems. One possibility is to treat the
system at a single-individual level, using ordinary differential equations
(ODEs) or delay differential equations (DDEs) to describe their motion in
space and time [22, 23, 24, 25]. Another possibility that is applicable with
sufficiently dense numbers of particles involves the use of continuum models
with averaged velocity and agent density fields that are governed by partial
differential equations (PDEs) [9, 26, 27, 11]. Moreoever, a number of
researchers have incorporated random noise effects into their models that are
able to produce transitions from one pattern to another [28, 29]. The study of
these systems has been enriched by tools from statistical physics [30] since
both first and second order phase transitions have been found in the formation
of coherent states [22, 31].
An important aspect of understanding self-organizing swarms patterns is that
of delay in the coupling between individual agents. Time delay appears in many
systems for several reasons: 1) finite time information transfer; 2) time
required to acquire measurement information; 3) computation time required for
generating the control instructions; and 4) actuation time required for the
instructions to be applied. In general, time delay reflects an important
property inherited in all swarms due to actuation, control, communication, and
computation [5, 32]. The occurrence of time delays in interacting particle
systems and in dynamical systems in general has been shown to have profound
dynamical consequences, such as destabilization and synchronization [33, 34].
Since time delays in the engineering of autonomous robot-systems are almost
unavoidable, incorporating them into the mathematical models is particularly
important. Initially, such studies focused on the case of one or a few
discrete time delays. More recently, however, the complex situation of several
and random time delays has been researched [35, 36, 37]. Another important
case is that of distributed time delays, when the dynamics of the system
depends on a continuous interval in its past instead of on a discrete instant
[38, 39, 40].
When examining self-organizing patterns in swarms, different attractors emerge
depending on initial conditions, and/or the the addition of external noise.
Co-existing bi- and multi-stable swarming patterns have been observed in a
multitude of models [41, 42, 43]. Because of the existence of co-existing
patterns, time delayed swarming systems display transitions between different
spatio-temporal patterns if there is an adequate balance between the strength
of the attractive coupling, the duration of time delay and the external noise
intensity, [29, 44, 45]. Often, mean-field approximations to swarm dynamics
cannot capture the possibility of different patterns for the same parameter
values [44]; this failure is a consequence of the mean field missing the
details of the particle distribution about the swarm center of mass. To
redress the bi-stable formation problem, we develop a new reduced model based
on a higher-order approximation that is able to predict bi-stable patterns in
globally coupled swarm models with delayed interactions.
## 2 The Reduced Swarm Model Derivation
We consider the dynamics of a two dimensional system of $N$ particles being
acted on by the influence of self-propulsion and mutual attraction. Let
$\mathbf{f}(\mathbf{r,\dot{r}})$, denote a self propelling force, and an
interaction potential function between particles be given by
$U(\mathbf{r}_{i}-\mathbf{r}_{j})$. In our description, the attraction between
particles does not occur instantaneously, but rather in a time delayed fashion
due to finite communication speeds and processing times. We describe the
general motion of the particles by the following dimensionless equations:
$\displaystyle\ddot{\mathbf{r}}_{i}(t)\mathbf{=\mathbf{f}}(\mathbf{\mathbf{r},\dot{r}}_{i})-\frac{a}{N}\sum_{j=1,j\neq
i}^{N}\nabla_{\mathbf{r}_{i}}U(\mathbf{r}_{i}(t)-\mathbf{r}_{j}(t-\tau))$ (1)
for $i=1,2\ldots,N$. Here, $\mathbf{r}_{i}(t)$ and $\dot{\mathbf{r}}_{i}(t)$
denote the two-dimensional position and velocity of particle $i$ at time $t$,
respectively. We assume
$\mathbf{f}(\mathbf{r},\dot{\mathbf{r}})=(1-\dot{\mathbf{r}}_{i}^{2})\dot{\mathbf{r}}_{i}$
to describe the self-propulsion of agent $i$, where
$\dot{\mathbf{r}}_{i}^{2}=\dot{\mathbf{r}}_{i}\cdot\dot{\mathbf{r}}_{i}$. We
denote the parameter $a$ as the coupling constant and measures the strength of
attraction between agents. At time $t$, agent $i$ is attracted to the position
of agent $j$ at the past time $t-\tau$. If we assume the form of our model is
based on the normal form for particles near a supercritical bifurcation
corresponding to the onset of coherent motion [46], then the leading term of
the potential function may be considered to be quadratic; i.e.,
$\nabla_{\mathbf{r}_{i}}U(\mathbf{r}_{i}(t)-\mathbf{r}_{j}(t-\tau))=\mathbf{r}_{i}(t)-\mathbf{r}_{j}(t-\tau)$.
Specific mathematical models of this kind have been extensively used to study
the dynamics of swarm patterns [47, 28, 48, 29, 46, 44]. Certainly, the choice
of potential function has a fundamental impact on the type of long-term
patterns that the system may acquire as well as on determining their
dimensionality and characteristic spatio-temporal length scales [49]; many of
these scales may be explicitly obtained for the patterns arising from Eqs. (1)
[44]. Numerous potential functions appropriate to different biological and
engineering situations have been carefully investigated [3, 2]; most of them
possess one or more minima where attraction and repulsion are balanced. Our
potential may be thought of as a first, quadratic approximation to the minima
of more complicated potential landscapes.
We obtain a reduced description for the swarm dynamics in Eqn. (1) by defining
the center of mass of the swarm (CM),
$\mathbf{R}(t)=\frac{1}{N}\sum_{i=1}^{N}\mathbf{r}_{i}(t)$, and the three
tensors
$\displaystyle\mathbf{C_{rr}}=\frac{1}{N}\sum_{i=1}^{N}\delta\mathbf{r}_{i}\delta\mathbf{r}_{i},\qquad\mathbf{C_{vv}}=\frac{1}{N}\sum_{i=1}^{N}\delta\dot{\mathbf{r}}_{i}\delta\dot{\mathbf{r}}_{i},$
(2)
$\displaystyle\mathbf{C_{rv}}=\frac{1}{N}\sum_{i=1}^{N}\delta\mathbf{r}_{i}\delta\dot{\mathbf{r}}_{i},$
where $\delta\mathbf{r}_{i}=\mathbf{r}_{i}-\mathbf{R}$. The above tensors
represent all of the second moments of the particles’ position and velocity
relative to the center of mass. Here, $\mathbf{a}\mathbf{b}$ is the exterior
product of the vectors $\mathbf{a}$ and $\mathbf{b}$ and has matrix components
$(\mathbf{a}\mathbf{b})_{ij}=a_{i}b_{j}$. Note that $\mathbf{C_{rr}}$ and
$\mathbf{C_{vv}}$ are symmetric tensors with non-negative diagonal elements,
whereas $\mathbf{C_{rv}}$ has neither of these properties.
The dynamical equation for the center of mass is obtained from the relation
$\ddot{\mathbf{R}}(t)=\frac{1}{N}\sum_{i=1}^{N}\ddot{\mathbf{r}}_{i}(t)$,
while the equations for the tensors are found by taking time derivatives of
$\delta\mathbf{r}_{i}$ and recalling $\sum_{i}\delta\mathbf{r}_{i}=0$. In our
derivation, we drop all possible third order moments (that take the form of
third order tensors) and justify it as follows. The swarm equations are
rotationally-invariant in space and so the time-asymptotic patterns that arise
tend to have particle position and velocity distributions that are symmetric
with respect to the CM. For large numbers of particles organized in such
symmetric patterns, these third order moments are composed of mutually
cancelling terms since they are of odd power in either in the position or the
velocity relative to the CM. Finally, we close the system of equations by
approximating fourth order moments of the form
$\frac{1}{N}\sum_{i=1}^{N}\delta\dot{\mathbf{r}}^{2}_{i}\mathbf{a}_{i}\mathbf{b}_{i}$
(where $\mathbf{a}_{i}$ and $\mathbf{b}_{i}$ are either
$\delta{\mathbf{r}}_{i}$ or $\delta\dot{\mathbf{r}}_{i}$) by
$\displaystyle\frac{1}{N}\sum_{i=1}^{N}\left\langle\delta\dot{\mathbf{r}}^{2}_{i}\right\rangle\mathbf{a}_{i}\mathbf{b}_{i}=\textrm{tr}(\mathbf{C_{vv}})\
\frac{1}{N}\sum_{i=1}^{N}\mathbf{a}_{i}\mathbf{b}_{i},$ (3)
and dropping all higher order moments. Here, $\textrm{tr}(\mathbf{C_{vv}})$
denotes the trace of $\mathbf{C_{vv}}$.
Our mean-field approximation including up to second moments (MF2M) finally
takes the form:
$\displaystyle\ddot{\mathbf{R}}=$ $\displaystyle
P(t)\dot{\mathbf{R}}-2\dot{\mathbf{R}}\cdot\mathbf{C_{vv}}-a\left(\mathbf{R}(t)-\mathbf{R}(t-\tau)\right),$
(4a) $\displaystyle\ddot{\mathbf{C}}_{\mathbf{rr}}=$ $\displaystyle
2\mathbf{C}_{\mathbf{vv}}-2\left(\dot{\mathbf{R}}\cdot{\mathbf{C}}_{\mathbf{vr}}\dot{\mathbf{R}}+\dot{\mathbf{R}}{\mathbf{C}}_{\mathbf{rv}}\cdot\dot{\mathbf{R}}\right)$
$\displaystyle\qquad+P(t)\dot{\mathbf{C}}_{\mathbf{rr}}-2a{\mathbf{C}}_{\mathbf{rr}},$
(4b) $\displaystyle\dot{\mathbf{C}}_{\mathbf{rv}}=$
$\displaystyle\mathbf{C}_{\mathbf{vv}}-2\dot{\mathbf{R}}\cdot{\mathbf{C}}_{\mathbf{vr}}\dot{\mathbf{R}}+P(t){\mathbf{C}}_{\mathbf{rv}}-a{\mathbf{C}}_{\mathbf{rr}},$
(4c) $\displaystyle\dot{\mathbf{C}}_{\mathbf{vv}}=$
$\displaystyle-2\left(\dot{\mathbf{R}}\cdot\mathbf{C_{vv}}\dot{\mathbf{R}}+\dot{\mathbf{R}}\mathbf{C_{vv}}\cdot\dot{\mathbf{R}}\right)+2P(t)\mathbf{C_{vv}}$
$\displaystyle\qquad-a\dot{\mathbf{C}}_{\mathbf{rr}},$ (4d)
where we let
$P(t)=\left(1-\dot{\mathbf{R}}^{2}-\textrm{tr}(\mathbf{C_{vv}})\right)$. Note
that no evolution equation is needed for the tensor $\mathbf{C}_{\mathbf{vr}}$
since $\mathbf{C}_{\mathbf{vr}}=\mathbf{C}_{\mathbf{rv}}^{\top}$.
Interestingly, no second moment tensor appears in a time-delayed form. This is
because the acceleration due to time delay that each particle and the center
of mass undergo is the same, up to $\mathcal{O}(1/N)$. Thus, relative to the
center of mass (as the second moment tensors are themselves measured) the
particles undergo no time-delayed acceleration, to the order mentioned. Also,
note that (_i_) Eqns. (4) reduce to the first-order mean-field approximation
of [44] when all second-moment tensors are neglected; and that (_ii_) the
linearization of Eqns. (4) about the trivial state decouples the CM and the
second moment tensors. The standard linear stability calculation shows that
second moment tensor equations render the trivial solution of Eqns. (4)
unstable for all parameter values. The instability of the trivial solution
agrees with what intuition tells us about Eqn. (1): the slightest difference
in position or velocity among the particles will accelerate them via self-
propulsive and attractive forces and make the CM and second-moment tensors
depart from the stationary solution. This effect is not captured by the simple
MF model since it does not account at all for how the particles are
distributed about the center of mass.
Figure 1: As in [44], Hopf (blue) and pitchfork (red) bifurcation curves in
$a$ and $\tau$ space that delimit different regions of long-term dynamical
behavior as prediceted by the first-order MF approximation: uniformly
travelling state (A), ring state (B) and rotating state (C). (Color online)
## 3 Model analysis and physical interpretation
The velocity second moment introduces two corrections into the propulsion of
the CM equation (4a) in a physically meaningful way. To see this, first note
that while the self-propulsive force of each individual particle always lies
along its velocity vector, the cumulative self-propulsion of all particles is
not necessarily directed along the CM velocity vector. The first tensor
correction accounts for how the particle dispersal slows down the CM
propulsion along its velocity and appears in the term
$P(t)\dot{\mathbf{R}}=\left(1-\dot{\mathbf{R}}^{2}-\textrm{tr}(\mathbf{C_{vv}})\right)\dot{\mathbf{R}}$.
For example, consider that all particles move coherently with the exact same
velocity vectors; since their self-propelling forces are also coherent the CM
self-propulsion term
$\left(1-\dot{\mathbf{R}}^{2}-\textrm{tr}(\mathbf{C_{vv}})\right)\dot{\mathbf{R}}$
is maximal, in the sense that $\textrm{tr}(\mathbf{C_{vv}})=0$. Otherwise,
when the particles are becoming dispersed $\textrm{tr}(\mathbf{C_{vv}})>0$,
their individual self-propulsion does not add up coherently and this makes the
self-propulsion of the CM weaker. The second tensor correction is the term
$-2\dot{\mathbf{R}}\cdot\mathbf{C_{vv}}$ and, since in contrast to the term
$P(t)\dot{\mathbf{R}}$ it is not necessarily directed along the velocity
vector $\dot{\mathbf{R}}$, it represents a correction for the fact that the CM
propulsion may have a component orthogonal to $\dot{\mathbf{R}}$ because of
the dispersal of particles.
An important consequence of (_i_) and (_ii_) is that all of the bifurcations
found previously for the MF approximation [44] are inherited by the MF2M
system in Eqns. (4). These bifurcation boundaries delimit the parameter
regions where the MF approximation predicts different spatio-temporal patterns
to be adopted in the long-time limit (Fig. 1). (A) A uniformly travelling
state is composed of particles collapsed together and moving at constant speed
in a given direction. (B) A ring state is formed by particles distributed
almost uniformly along a circle, some of them moving clockwise and others
counterclockwise, while the center of mass is at rest at the center (see Fig.
4a below). (C) A rotating state, in which all particles collapse to a point
and move in a circular orbit (see Fig. 4b).
The spatio-temporal patterns (A), (B) and (C) are captured by our MF2M
approximation and manifest themselves as follows. The uniformly
travelling/rotating states have trivial components for the second moment
tensors (indicating the collapse of all particles to the CM) but uniform
motion/periodic oscillations for the position of the CM. In contrast, the ring
state has a stationary position for the CM but periodic oscillations for all
second moment tensor components. The periodic oscillations of the tensor
components in the ring state are due to the fact that particles are not
distributed quite uniformly along the ring in either position or velocity
space and their spread about the CM (second moment tensors) in both spaces has
periodic variations.
Figure 2: Period of sustained oscillations for the first component of the CM
vector (left column) and the tensor component $C_{\mathbf{vv}}(1,1)$ (right
column) for different values of the time delay $\tau$ and for $a=2$. Blue
circles correspond to the full swarm described by Eqn. (1) and red crosses to
our MF2M approximation. Initial conditions favor convergence to the
rotating/ring state along the top/bottom row. The swarm converges to the ring
state for $\tau\lesssim 1.5$ ($\tau\lesssim 2.1$) along the top (bottom) row
and to the rotating state for higher $\tau$’s. In the ring state, the period
of oscillation of $C_{\mathbf{vv}}(1,1)$ is approximately $\pi/\sqrt{a}=2.22$,
as shown analytically. The agreement with other tensor components is similar.
See text for more details. (Color online)
We employ numerical simulations using a fourth order Runge-Kutta algorithm
with quadratic interpolation for the delayed terms to see how our MS2M system
captures the convergence to different spatio-temporal patterns. The coupling
parameter is fixed at $a=2$ and the time delay is varied so as to cross the
boundary between regions B and C of Fig. 1. We consider two types of initial
conditions over the interval $[-\tau,0]$: in the first, the particles are
distributed at random in a box with sides of length 0.05 and velocities
uniformly distributed in [0.5 0.55] for the $x$ component and in [-0.025
0.025] for the $y$ component. In the second initial condition we distribute
$N=300$ particles at random in the unit box with velocity components randomly
distributed between -0.5 and 0.5. The first initial condition favors
convergence to the rotating state while the second one favors convergence to
the ring state. We note that although the initial conditions for the position
and velocity vectors appear physically inconsistent, this is of no consequence
since no time-delayed velocity terms appear in the equations.
Starting from $\tau=0.5$, the full swarm and our MF2M systems first converge
asymptotically to the ring state (Fig. 2). As the time delay increases, the
convergence is to the rotating state instead. However, we find the hysterisis
loop characteristic of bi-stable behavior: the transition occurs at different
values of $\tau$ for the two different initial conditions. Remarkably, our
MF2M system is not only able to accurately predict the periods of oscillation
for the ring/rotating states of the full swarm equations, but it also
identifies the value of the time delay at which the pattern transition occurs
for each of the two initial conditions.
Figure 3: Oscillation amplitude of the first component of the CM vector (left
column) and the tensor component $C_{\mathbf{vv}}(1,1)$ (right column), after
the decay of transients, for different values of the time delay $\tau$ and for
$a=2$. Blue circles correspond to the full swarm described by Eqn. (1) and red
crosses to our MF2M approximation. Initial conditions favor convergence to the
rotating/ring state along the top/bottom row. The swarm converges to the ring
state for $\tau\lesssim 1.5$ ($\tau\lesssim 2.1$) along the top (bottom) row
and to the rotating state for higher $\tau$’s. The agreement with other tensor
components is similar. (Color online)
As noted above, the second moment tensors undergo periodic oscillations when
the particles adopt the ring state. The period of these tensor oscillations
may be approximated directly from Eqns. (4). In this spatio-temporal pattern,
the CM is fixed and individual particles move at unit speed [44]. Substituting
$\mathbf{R}=\dot{\mathbf{R}}=0$ and $\textrm{tr}(\mathbf{C_{vv}})=1$ in Eqns.
(4) yields a linear system with oscillatory solutions with period
$\pi/\sqrt{a}$. On the other hand, the period of oscillation of the center of
mass position in the rotating state is determined by a complicated non-linear
equation that may be obtained by resorting to polar coordinates and agrees
perfectly with numerical simulations [44]. Thus, while the simple MF
approximation fails to capture the bi-stable behavior displayed by the full
swarm model, our MF2M approximation is able to do so.
Figure 4: Top panel: hysterisis loop of the CM speed for the full model (solid
blue) and our MF2M (dashed red) approximation, obtained with a time varying
time delay (loop is traversed counter-clockwise in time). A zero/non-zero CM
speed indicates convergence to the ring/rotating state. Transient dynamics at
the time delay values where pattern switching occurs ($\tau=1.1$, 2.4) have
been removed. Bottom panels: time snapshots of the spatial distribution of the
swarm particles at the points labeled (a) and (b) of the top panel. In the
rotating state of panel (b) all particles have collapsed to a point and move
with the same velocity (blue arrow) around a circle; a portion of their time
track is also shown (red line). (Color online).
An important signature of the spatio-temporal patterns of our system is the
amplitude of oscillations, which we also extract from the results of our
numerical simulations (Fig. 3). The hysterisis of the system becomes equally
evident in these amplitude vs. time delay plots. For the two initial
conditions described before, our MF2M model is able to capture the oscillation
amplitude of the CM (Fig. 3, left column) with accuracy as well as the time
delay value at which the long-time convergence switches from the ring to the
rotating state. However, the amplitude of oscillation of the tensor component
$C_{\mathbf{vv}xx}$ is not captured as well, particularly for time delay
values near the MF bifurcation (Fig. 3, right column). This departure is due
to the neglected higher-order moments and finite-particle effects.
An additional way to visualize the bistable nature of the ring and rotating
state attractors is by forcing the system to undergo hysterisis via a time-
dependent time delay. To this end, we start both the full swarm system and our
MF2M approximation with a time delay of $\tau=1.1$, increase it slowly up to
$\tau=2.4$ and then bring it back down to its starting value. Using the speed
of the CM as a proxy for the state of the system, we see the swarm converge to
the ring state (CM speed of zero) up until the time delay reaches 2.4; at this
point, the swarm switches to the rotating state (non-zero CM speed) and
persists in this state until the time delay is back to 1.1 (Fig. 4, top
panel). Figs. 4(a) and (b) show snapshots of the swarm in different states,
taken at the points indicated in the top panel.
## 4 Discussion and conlusions
A more realistic version of the model studied here should modify our choice of
potential function to include the local, hard-body repulsion among
individuals. Our previous work shows that the patterns and transitions
discussed here do not fundamentally change with the addition of small, local
repulsive forces between particles. Stronger repulsion certainly can
destabilize the coherent structures. Crucially, we note that since repulsion
causes the particles to spread out in space, their distribution about the
center of mass becomes even more important for determining their group
dynamics. Although a systematic study is beyond the scope of this work, we
expect our extended MF2M model to be highly useful in frameworks where
repulsion between particles cannot be neglected.
We also expect our multi-moment approach to extend the MF2M model framework to
situations in which the time delay for inter-particle interactions is
distributed. This situation is of interest since, analogously to inter-
particle repulsion, randomly distributed time delays cause the individual
agents to occupy larger portions of space, making their second moments non-
negligible [45]. Conversely, in some realistic applications the particle
distribution about the center of mass can have an effect on the magnitude of
the time delay. In our MF2M framework, this would make the time delay be
dependent on the second moment tensors.
To summarize, we derived a new MF approximation of a general model swarming
system in order to account for higher order moments about the CM. Notably, our
extended MF2M model is able to account for the bi-stability of spatio-temporal
patterns that are displayed by the original swarming particles. This is in
sharp contrast to the MF approximation for this system, which cannot capture
this bi-stability and other complex behavior of the time delayed swarming
system that we studied. With the inclusion of higher order moments, it is
clear to what extent bi-stability depends on the particle distribution of the
swarm about its center of mass. Adding the additional physics to the model,
although higher dimensional than just a mean field, should allow more general
low dimensional approaches to more accurately predict the structure of bi-
stability in large population swarms.
###### Acknowledgements.
The authors gratefully acknowledge the Office of Naval Research for their
support through ONR contract no. N0001412WX20083 and the Naval Research
Laboratory 6.1 program contract no. N0001412WX30002. LMR and IBS are supported
by Award Number R01GM090204 from the National Institute Of General Medical
Sciences. The content is solely the responsibility of the authors and does not
necessarily represent the official views of the National Institute Of General
Medical Sciences or the National Institutes of Health.
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|
arxiv-papers
| 2013-09-19T18:05:19 |
2024-09-04T02:49:51.189858
|
{
"license": "Public Domain",
"authors": "Luis Mier-y-Teran-Romero and Ira B. Schwartz",
"submitter": "Ira Schwartz",
"url": "https://arxiv.org/abs/1309.5064"
}
|
1309.5189
|
# On the similarity of Tensors111P. Yuan’s research is supported by the NSF of
China (Grant No. 11271142) and the Guangdong Provincial Natural Science
Foundation(Grant No. S2012010009942), L. You’s research is supported by the
Zhujiang Technology New Star Foundation of Guangzhou (Grant No. 2011J2200090)
and Program on International Cooperation and Innovation, Department of
Education, Guangdong Province (Grant No.2012gjhz0007).
Pingzhi Yuan222Email address: [email protected]., Lihua You333Corresponding
author: [email protected].
(School of Mathematical Sciences, South China Normal University,
Guangzhou, 510631, P.R. China
)
###### Abstract
In this paper, we characterize all similarity relations when $m\geq 3$, obtain
some interesting properties which are different from the case $m=2$, and show
that the results of matrices about the Jordan canonical form cannot be
extended to tensors. AMS classification: 15A18; 15A69 Keywords: tensor;
product; similarity; hypergraph.
## 1 Introduction
Since the work of Qi [5] and Lim [3], the study of tensors and the spectra of
tensors (and hypergraphs) and their various applications have attracted much
attention and interest.
An order $m$ dimension $n$ tensor $\mathbb{A}=(a_{i_{1}i_{2}\ldots
i_{m}})_{1\leq i_{j}\leq n\hskip 5.69046pt(j=1,\ldots,m)}$ over the complex
field $\mathbb{C}$ is a multidimensional array with all entries
$a_{i_{1}i_{2}\ldots
i_{m}}\in\mathbb{C}\,(i_{1},\ldots,i_{m}\in[n]=\\{1,\ldots,n\\})$. The
majorization matrix $M(\mathbb{A})$ of the tensor $\mathbb{A}$ is defined as
$(M(\mathbb{A}))_{ij}=a_{ij\ldots j},(i,j\in[1,n])$ by Pearson [4]. The unit
tensor of order $m$ and dimension $n$ is the tensor
$\mathbb{I}=(\delta_{i_{1},i_{2},\ldots,i_{m}})$ with
$\delta_{i_{1},i_{2},\ldots,i_{m}}=\left\\{\begin{array}[]{cc}1,&{\rm if\hskip
5.69046pt}i_{1}=i_{2}=\ldots=i_{m};\\\ 0,&{\rm otherwise}.\end{array}\right.$
Let $\mathbb{A}$ (and $\mathbb{B}$) be an order $m\geq 2$ (and $k\geq 1$),
dimension $n$ tensor, respectively. Recently, Shao [6] defined a general
product $\mathbb{A}\mathbb{B}$ to be the following tensor $\mathbb{D}$ of
order $(m-1)(k-1)+1$ and dimension $n$:
$d_{i\alpha_{1}\ldots\alpha_{m-1}}=\sum\limits_{i_{2},\ldots,i_{m}=1}^{n}a_{ii_{2}\ldots
i_{m}}b_{i_{2}\alpha_{1}}\ldots
b_{i_{m}\alpha_{m-1}}\quad(i\in[n],\,\alpha_{1},\ldots,\alpha_{m-1}\in[n]^{k-1}).$
The tensor product possesses a very useful property: the associative law ([6],
Theorem 1.1). With the general product, the following definition of the
similarity relation of two tensors was proposed by Shao [6].
###### Definition 1.1.
([6], Definition 2.3) Let $\mathbb{A}$ and $\mathbb{B}$ be two order $m$
dimension $n$ tensors. Suppose that there exist two matrices $P$ and $Q$ of
order $n$ with $P\mathbb{I}Q=\mathbb{I}$ such that $\mathbb{B}=P\mathbb{A}Q$,
then we say that the two tensors are similar.
It is easy to see that the similarity relation is an equivalent relation, and
similar tensors have the same characteristic polynomials, and thus they have
the same spectrum (as a multiset). For example, the permutation similarity and
the diagonal similarity (see also [6, 7, 8]) are two special kinds of the
similarity of tensors.
###### Definition 1.2.
([6]) Let $\mathbb{A}$ and $\mathbb{B}$ be two order $m$ dimension $n$
tensors. We say that $\mathbb{A}$ and $\mathbb{B}$ are permutational similar,
if there exists some permutation matrix $P$ of order $n$ such that
$\mathbb{B}=P\mathbb{A}P^{T}$, where $\sigma\in S_{n}$ is a permutation on the
set $[n]$ and $P=P_{\sigma}=(p_{ij})$ is the corresponding permutation matrix
of $\sigma$ with $p_{ij}=1\Leftrightarrow j=\sigma(i)$.
###### Definition 1.3.
([6], Definition 2.4) Let $\mathbb{A}$ and $\mathbb{B}$ be two order $m$
dimension $n$ tensors. We say that $\mathbb{A}$ and $\mathbb{B}$ are diagonal
similar, if there exists some invertible diagonal matrix $D$ of order $n$ such
that $\mathbb{B}=D^{-(m-1)}\mathbb{A}Q$.
About the matrices $P$ and $Q$ in Definition 1.1, [6] showed the following
proposition.
###### Proposition 1.4.
([6], Remark 2.1) If $P$ and $Q$ are two matrices of order $n$ with
$P\mathbb{I}Q=\mathbb{I}$, where $\mathbb{I}$ is the order $m$ dimension $n$
unit tensor, then both $P$ and $Q$ are invertible matrices.
In this paper, we characterize all similarity relations when $m\geq 3$ in
Section 2, and in Section 3, we obtain some interesting properties which are
different from the case $m=2$, and show that the results of matrices about the
Jordan canonical form cannot be extended to tensors.
## 2 Main results
In this section, we will prove Theorem 2.3. We first prove the following two
lemmas.
###### Lemma 2.1.
Let $P=(p_{ij})$ and $Q$ be two matrices of order $n$ with
$P\mathbb{I}Q=\mathbb{I}$, $\mathbb{I}$ the order $m$ dimension $n$ unit
tensor. Take $\mathbb{I}Q=\mathbb{A}=(a_{i_{1}i_{2}\ldots i_{m}})$, then we
have
(i) $a_{i_{1}i_{2}\ldots i_{m}}=0$ when
$(i_{2},i_{3},\ldots,i_{m})\neq(i_{2},i_{2},\ldots,i_{2})$;
(ii) $PM(\mathbb{A})=I$, where $M(\mathbb{A})$ is the majorization matrix of
$\mathbb{A}$ and $I$ is the unit matrix of order $n$.
###### Proof.
For the proof of (i), we let $\alpha=i_{2}i_{3}\ldots i_{m}\neq
i_{2}i_{2}\ldots i_{2}$. By $P\mathbb{A}=\mathbb{I}$, we have
$\delta_{i\alpha}=\sum_{j=1}^{n}p_{ij}a_{j\alpha},\,i=1,2,\ldots,n.$
Since $\alpha\neq i_{2}i_{2}\ldots i_{2}$, then $\delta_{i\alpha}=0$, it
follows that
$P(a_{1\alpha},a_{2\alpha},\ldots,a_{n\alpha})^{T}=(0,0,\ldots,0)^{T}.$
Observe that $P$ is an invertible matrix by Proposition 1.4, thus
$a_{j\alpha}=0,j=1,2,\ldots,n$, and the conclusion of (i) follows.
For the proof of (ii), by $P\mathbb{A}=\mathbb{I}$, we have
$\delta_{ij\ldots j}=\sum_{u=1}^{n}p_{iu}a_{uj\ldots
j}=\sum_{u=1}^{n}p_{iu}(M(\mathbb{A}))_{uj}=(PM(\mathbb{A}))_{ij}.$
Hence $PM(\mathbb{A})=I$, this completes the proof of (ii).∎
###### Lemma 2.2.
Let $\mathbb{I}$ be the order $m\geq 3$ dimension $n$ unit tensor. Suppose
that $Q=(q_{ij})$ is a matrix of order $n$ such that
$(\mathbb{I}Q)_{i\alpha}=0$
for all $i\in[n]$ and all $\alpha\neq j\ldots j,j\in[n]$, then there is at
most one nonzero element in every row of $Q$.
###### Proof.
Since
$(\mathbb{I}Q)_{i_{1}i_{2}\ldots
i_{m}}=\sum_{j_{2},\ldots,j_{m}=1}^{n}\delta_{i_{1}j_{2}\ldots
j_{m}}q_{j_{2}i_{2}}\ldots q_{j_{m}i_{m}}=q_{i_{1}i_{2}}q_{i_{1}i_{3}}\ldots
q_{i_{1}i_{m}},$
by the assumption, for every $\alpha=i_{2}i_{3}\ldots i_{m}\neq
i_{2}i_{2}\ldots i_{2}$, we get
$q_{i_{1}i_{2}}q_{i_{1}i_{3}}\ldots q_{i_{1}i_{m}}=0.$ (2.1)
If $q_{i_{1}t}\neq 0$ for some $t\in[n]\setminus\\{i_{2}\\}$, then we take
$i_{3}=\ldots=i_{m}=t$. By (2.1) we have $q_{i_{1}i_{2}}q_{i_{1}t}^{m-2}=0$,
hence $q_{i_{1}i_{2}}=0$. Since the choice of $i_{2}$ is arbitrary, then we
have proved that for any $i\in[n]$, there is at most one nonzero element in
$\\{q_{i1},q_{i2},\ldots,q_{in}\\}$. This completes the proof of the lemma.∎
###### Theorem 2.3.
Let $\mathbb{I}$ be the unit tensor of order $m\geq 3$ and dimension $n\geq
2$. Suppose that $P$ and $Q$ are two matrices of order $n$ with
$P\mathbb{I}Q=\mathbb{I}$, then there exist a permutation matrix $R_{\sigma}$
and an invertible diagonal matrix $D$, such that $Q=DR_{\sigma}$ and
$P=R_{\sigma}^{T}D^{1-m}$, where $\sigma\in S_{n}$ is a permutation on the set
$[n]$.
###### Proof.
Since $\mathbb{I}=P\mathbb{I}Q=P(\mathbb{I}Q)$, combining Lemmas 2.1 and 2.2,
we obtain that there is at most one nonzero element in every row of $Q$. Note
that $Q$ is an invertible matrix by Proposition 1.4, hence there is precisely
one nonzero element, say, $q_{i\sigma(i)}\neq 0,1\leq i\leq n$ in every row of
$Q$. Obviously, $\sigma(1),\sigma(2),\ldots,\sigma(n)$ is a permutation of
$\\{1,2,\ldots,n\\}$, hence
$Q=DR_{\sigma}$
is the product of an invertible diagonal matrix
$D=diag(d_{11},\ldots,d_{nn})$, where $d_{ii}=q_{i\sigma(i)}$ and a
permutation matrix $R_{\sigma}=(r_{ij})$, where $r_{ij}=1$ if $j=\sigma(i)$
and $r_{ij}=0$ otherwise.
Now by the general tensor product, we get
$(M(\mathbb{I}Q))_{ij}=(\mathbb{I}Q)_{ij\ldots
j}=q_{ij}^{m-1}=\left\\{\begin{array}[]{cc}q_{i\sigma(i)}^{m-1}&{\rm if\hskip
5.69046pt}j=\sigma(i);\\\ 0&{\rm otherwise}.\end{array}\right.$
Hence $M(\mathbb{I}Q)$ is the product of the invertible diagonal matrix
$D^{m-1}=diag(d_{11}^{m-1},\ldots,d_{nn}^{m-1})$ and the permutation matrix
$R_{\sigma}=(r_{ij})$. By Lemma 2.1 (ii), we have $PM(\mathbb{I}Q)=I$, it
follows that $P=R_{\sigma}^{T}D^{1-m}$, is the product of the permutation
matrix $R_{\sigma}^{T}$ and the invertible diagonal matrix
$D^{1-m}=diag(d_{11}^{1-m},\ldots,d_{nn}^{1-m})$. ∎
###### Remark 2.4.
By the proof of Theorem 2.3, we see that $Q$ and $P$ are closely related with
$QP=D^{2-m}$ and $PQ=R_{\sigma}^{T}D^{2-m}R_{\sigma}$ are invertible diagonal
matrices. Note that when $m=2$, $P,Q$ are invertible matrices and $PQ=QP=I$.
It implies that the case of $m\geq 3$ is completely different from the case
$m=2$.
Let $\mathbf{P}_{n}$ be the set of all permutation matrices of order $n$,
$\mathbb{P}_{n}$ be the set of all matrices which have the same zero patterns
with some permutation matrix of order $n$. Clearly,
$\mathbf{P}_{n}\subseteq\mathbb{P}_{n}$. For example,
$S=\left(\begin{array}[]{cc}0&1\\\ 1&0\\\
\end{array}\right)\in\mathbf{P}_{2},T=\left(\begin{array}[]{cc}0&2\\\ -3&0\\\
\end{array}\right)\in\mathbb{P}_{2}$, where $S$ and $T$ have the same zero
pattern.
###### Remark 2.5.
Let $\mathbb{I}$ be the unit tensor of order $m\geq 3$ and dimension $n\geq
2$. Suppose that $P$ and $Q$ are two matrices of order $n$ with
$P\mathbb{I}Q=\mathbb{I}$, then $P,Q\in\mathbb{P}_{n}$.
###### Remark 2.6.
By Definition 1.1, we know that Theorem 2.3 gives a characterization for the
similarities of tensors with order $m\geq 3$ dimension $n$, i.e., as for the
similarity of tensors, we need only consider the permutation similarity, the
diagonal similarity and their compositions.
###### Theorem 2.7.
Let $\mathbb{A}$ and $\mathbb{B}$ be two order $m\geq 3$ dimension $n$
tensors. If the tensors $\mathbb{A}$ and $\mathbb{B}$ are similar, then there
exists a tensor $\mathbb{C}$ such that $\mathbb{A}$ and $\mathbb{C}$ are
diagonal similar, and $\mathbb{B}$ and $\mathbb{C}$ are permutational similar.
###### Proof.
By Definition 1.1 and Theoren 2.3, there exist a permutation matrix $R$ and an
invertible diagonal matrix $D$ of order $n$ such that $Q=DR$, $P=R^{T}D^{1-m}$
and $\mathbb{B}=P\mathbb{A}Q$.
Take $\mathbb{C}=D^{1-m}\mathbb{A}D$, then
$\mathbb{B}=R^{T}\mathbb{C}R=R^{T}\mathbb{C}(R^{T})^{T}$. By Definitions 1.2
and 1.3, the results hold. ∎
## 3 Some applications
Let $Z(\mathbb{A})$ be the tensor obtained by replacing all the nonzero
entries of $\mathbb{A}$ by one. Then $Z(\mathbb{A})$ is called the zero-
nonzero pattern (or simply the zero pattern) of $\mathbb{A}$. Let $a$ be a
complex number, we define $Z(a)=1$ if $a\not=0$ and $Z(a)=0$ if $a=0$.
###### Lemma 3.1.
Let $\mathbb{A}=(a_{i_{1}i_{2}\ldots i_{m}})_{1\leq i_{j}\leq n\hskip
3.41418pt(j=1,\ldots,m)}$ and $\mathbb{B}=(b_{i_{1}i_{2}\ldots i_{m}})_{1\leq
i_{j}\leq n\hskip 2.84544pt(j=1,\ldots,m)}$ be two order $m\geq 3$ dimension
$n$ tensors. If the tensors $\mathbb{A}$ and $\mathbb{B}$ are diagonal
similar, then $Z(\mathbb{A})$=$Z(\mathbb{B})$.
###### Proof.
By Definition 1.3, there exists an invertible diagonal matrix
$D=diag(d_{11},\ldots,d_{nn})$ of order $n$ such that
$\mathbb{B}=D^{1-m}\mathbb{A}D$. Then
$b_{i_{1}i_{2}\ldots i_{m}}=(D^{1-m}\mathbb{A}D)_{i_{1}i_{2}\ldots i_{m}}$
$=\sum\limits_{j_{1},\ldots,j_{m}=1}^{n}(D^{1-m})_{i_{1}j_{1}}a_{j_{1}j_{2}\ldots
j_{m}}d_{j_{2}i_{2}}\ldots d_{j_{m}i_{m}}$
$=a_{i_{1}i_{2}\ldots i_{m}}d^{1-m}_{i_{1}i_{1}}d_{i_{2}i_{2}}\ldots
d_{i_{m}i_{m}}.$
Therefore $b_{i_{1}i_{2}\ldots i_{m}}\not=0\Leftrightarrow a_{i_{1}i_{2}\ldots
i_{m}}\not=0$, and thus $Z(\mathbb{A})$=$Z(\mathbb{B})$. ∎
###### Theorem 3.2.
Let $\mathbb{A}$ and $\mathbb{B}$ be two order $m\geq 3$ dimension $n$
tensors. If the tensors $\mathbb{A}$ and $\mathbb{B}$ are similar, then
$Z(\mathbb{A})$ and $Z(\mathbb{B})$ are permutational similar.
###### Proof.
By Theorem 2.7, there exists a tensor $\mathbb{C}$ such that $\mathbb{A}$ and
$\mathbb{C}$ are diagonal similar, and $\mathbb{B}$ and $\mathbb{C}$ are
permutational similar. It is easy that $Z(\mathbb{A})=Z(\mathbb{C})$ by Lemma
3.1. Now we show that $Z(\mathbb{B})$ and $Z(\mathbb{C})$ are permutational
similar.
Let $R$ be a permutation matrix of order $n$ such that
$\mathbb{B}=R\mathbb{C}R^{T}$. Then
$(Z(\mathbb{B}))_{i_{1}i_{2}\ldots
i_{m}}=(Z(R\mathbb{C}R^{T}))_{i_{1}i_{2}\ldots i_{m}}$
$=Z(\sum\limits_{j_{1},\ldots,j_{m}=1}^{n}r_{i_{1}j_{1}}c_{j_{1}j_{2}\ldots
j_{m}}(R^{T})_{j_{2}i_{2}}\ldots(R^{T})_{j_{m}i_{m}})$
$=Z(\sum\limits_{j_{1},\ldots,j_{m}=1}^{n}r_{i_{1}j_{1}}c_{j_{1}j_{2}\ldots
j_{m}}r_{i_{2}j_{2}}\ldots r_{i_{m}j_{m}})$
$=Z(c_{\sigma(i_{1})\sigma(i_{2})\ldots\sigma(i_{m})})$
$=\sum\limits_{j_{1},\ldots,j_{m}=1}^{n}r_{i_{1}j_{1}}Z(c_{j_{1}j_{2}\ldots
j_{m}})r_{i_{2}j_{2}}\ldots r_{i_{m}j_{m}}.$
So $Z(\mathbb{B})=RZ(\mathbb{C})R^{T}$, and thus $Z(\mathbb{A})$ and
$Z(\mathbb{B})$ are permutational similar. ∎
###### Corollary 3.3.
Let $\mathbb{A}$ and $\mathbb{B}$ be two order $m\geq 3$ dimension $n$
tensors, $N(\mathbb{A})$ the number of the nonzero entries of tensor
$\mathbb{A}$. If the tensors $\mathbb{A}$ and $\mathbb{B}$ are similar, then
$N(\mathbb{A})=N(\mathbb{B})$.
###### Proof.
By Theorem 3.2, there exists a permutation matrix $R=R_{\sigma}=(r_{ij})$ of
order $n$ such that $Z(\mathbb{B})=RZ(\mathbb{A})R^{T}$, where $\sigma\in
S_{n}$ is a permutation on the set $[n]$ and $r_{ij}=1\Leftrightarrow
j=\sigma(i)$. Then similar to the proof of Theorem 3.2, we have
$(Z(\mathbb{B}))_{i_{1}i_{2}\ldots
i_{m}}=(RZ(\mathbb{A})R^{T})_{i_{1}i_{2}\ldots i_{m}}$
$=\sum\limits_{j_{1},\ldots,j_{m}=1}^{n}r_{i_{1}j_{1}}(Z(\mathbb{A}))_{j_{1}j_{2}\ldots
j_{m}}(R^{T})_{j_{2}i_{2}}\ldots(R^{T})_{j_{m}i_{m}}$
$=\sum\limits_{j_{1},\ldots,j_{m}=1}^{n}r_{i_{1}j_{1}}(Z(\mathbb{A}))_{j_{1}j_{2}\ldots
j_{m}}r_{i_{2}j_{2}}\ldots r_{i_{m}j_{m}}$
$=(Z(\mathbb{A}))_{\sigma(i_{1})\sigma(i_{2})\ldots\sigma(i_{m})}.$
Therefore $(Z(\mathbb{B}))_{i_{1}i_{2}\ldots
i_{m}}\not=0\Leftrightarrow(Z(\mathbb{A}))_{\sigma(i_{1})\sigma(i_{2})\ldots\sigma(i_{m})}\not=0$,
and thus $N(\mathbb{A})$=$N(\mathbb{B})$. ∎
###### Remark 3.4.
Note that the result of Corollary 3.3 does not hold when $m=2$. For example,
let $P=\left(\begin{array}[]{cc}1&1\\\ 1&0\\\ \end{array}\right),$
$Q=\left(\begin{array}[]{cc}1&-1\\\ 1&0\\\ \end{array}\right),$
$A=\left(\begin{array}[]{cc}0&0\\\ 1&0\\\ \end{array}\right),$
$B=PAQ=\left(\begin{array}[]{cc}1&-1\\\ 1&-1\\\ \end{array}\right).$ It is
clear $PIQ=QIP=I$, then the matrix $A$ and $B$ are similar, but
$N(A)=1\not=N(B)=4.$
###### Definition 3.5.
([1]) Let $\mathbb{A}=(a_{i_{1}i_{2}\ldots i_{m}})_{1\leq i_{j}\leq n\hskip
3.41418pt(j=1,\ldots,m)}$ be an order $m\geq 3$ dimension $n$ tensor. If
$a_{i_{1}i_{2}\ldots i_{m}}\equiv 0$ whenever
$\min\\{i_{2},\ldots,i_{m}\\}<i_{1}$, then $\mathbb{A}$ is called an upper
triangular tensor. If $a_{i_{1}i_{2}\ldots i_{m}}\equiv 0$ whenever
$\max\\{i_{2},\ldots,i_{m}\\}>i_{1}$, then $\mathbb{A}$ is called a low
triangular tensor. If $\mathbb{A}$ is either upper or low triangular tensor,
then $\mathbb{A}$ is called a triangular tensor. If $a_{i_{1}i_{2}\ldots
i_{m}}\equiv 0$ whenever $i_{1}i_{2}\ldots i_{m}\not=i_{1}i_{1}\ldots i_{1}$,
then $\mathbb{A}$ is called a diagonal tensor.
Clearly, a triangular tensor is both an upper and a low triangular tensor.
###### Corollary 3.6.
Let $\mathbb{A}$ and $\mathbb{B}$ be two order $m\geq 3$ dimension $n$
tensors. Suppose that $\mathbb{A}$ and $\mathbb{B}$ are similar, and
$\mathbb{A}$ is a diagonal tensor, then $\mathbb{B}$ is a diagonal tensor.
###### Proof.
We only need to show $(Z(\mathbb{B}))_{i_{1}i_{2}\ldots i_{m}}=0$ (and thus
$b_{i_{1}i_{2}\ldots i_{m}}=0)$ for any $i_{1}i_{2}\ldots
i_{m}\not=i_{1}i_{1}\ldots i_{1}$ where $i_{1}\in[n]$.
By Theorem 3.2, there exists a permutation matrix $R=R_{\sigma}=(r_{ij})$ of
order $n$ such that $Z(\mathbb{B})=RZ(\mathbb{A})R^{T}$, where $\sigma\in
S_{n}$ is a permutation on the set $[n]$ and $r_{ij}=1\Leftrightarrow
j=\sigma(i)$. Then by the proof of Corollary 3.3, we have
$(Z(\mathbb{B}))_{i_{1}i_{2}\ldots
i_{m}}\not=0\Leftrightarrow(Z(\mathbb{A}))_{\sigma(i_{1})\sigma(i_{2})\ldots\sigma(i_{m})}\not=0$.
Note that $i_{1}i_{2}\ldots i_{m}\not=i_{1}i_{1}\ldots i_{1}$ if and only if
$\sigma(i_{1})\sigma(i_{2})\ldots\sigma(i_{m})\not=\sigma(i_{1})\sigma(i_{1})\ldots\sigma(i_{1})$,
therefore $(Z(\mathbb{B}))_{i_{1}i_{2}\ldots i_{m}}=0$ when $i_{1}i_{2}\ldots
i_{m}\not=i_{1}i_{1}\ldots i_{1}$ since $\mathbb{A}$ is a diagonal tensor. ∎
###### Remark 3.7.
Note that the result of Corollary 3.6 does not hold when $m=2$. For example,
all real symmetric matrices are similar to diagonal matrices, but symmetric
matrices are not diagonal matrices in general. Another example, if a real
matrix $A$ of order $n$ has $n$ distinct eigenvalues, then $A$ is similar to
diagonal matrices, but $A$ is not necessarily a diagonal matrix.
###### Definition 3.8.
A Jordan block $J_{k}(\lambda)$ is an upper triangular matrix of order $k$ of
the form
$P=\left(\begin{array}[]{ccccc}\lambda&1&&&\\\ &\lambda&1&&\\\
&&\ddots&\ddots&\\\ &&&\ddots&1\\\ &&&&\lambda\end{array}\right).$
There are $k-1$ terms 1 in the superdiagonal; the scalar $\lambda$ appears $k$
times on the main diagonal. All other entries are zero, and
$J_{1}(\lambda)=[\lambda]$.
###### Theorem 3.9.
([2], Jordan canonical form theorem) Let $A$ be a given complex matrix of
order $n$. There is a nonsingular matrix $S$ of order $n$ such that
$A=SJS^{-1}$, where
$J=\left(\begin{array}[]{ccccc}J_{n_{1}}(\lambda_{1})&&&0\\\
&J_{n_{2}}(\lambda_{2})&&\\\ &&\ddots&\\\
0&&&J_{n_{k}}(\lambda_{k})\end{array}\right)$
is a Jordon matrix with $n_{1}+n_{2}+\ldots+n_{k}=n$, and $J$ is unique up to
permutations of the diagonal Jordan blocks, called the Jordan canonical form
of $A$. The eigenvalues $\lambda_{i},i=1,\ldots,k$ are not necessarily
distinct. If $A$ is a real matrix with only real eigenvalues, then the
similarity matrix $S$ can be taken to be real.
Clearly, the Jordon matrix $J$ is an upper triangular matrix.
###### Remark 3.10.
The results of matrices about the Jordan canonical form cannot be extended to
tensors since not all tensors are similar to some upper triangular tensor by
Corollary 3.3.
## References
* [1] S.L. Hu, Z.H. Huang, C. Ling, L.Q. Qi, On determinants and eigenvalue theory of tensors, Journal of Symbolic Computation, 50(2013), 508-531.
* [2] R.A. Horn, C.R. Johnson, Matrix Analysis, Posts and Telecom Press, 2005.
* [3] L.H. Lim, Singular values and eigenvalues of tensors, a rational approach, in proceedings 1st IEEE international workshop on computational advances of adaptive processing (2005), 129-132.
* [4] K. Pearson, Essentially positive tensors, Int. J. Algebra, 4(2010),421-427.
* [5] L. Qi, Eigenvalues of a real supersymmetric tensor, Symbolic Comput. 40(2005), 1302-1324.
* [6] J.Y Shao, A general product of tensors with applications, Linear Algebra and its Appl. 439(2013), 2350-2366.
* [7] Y Yang and Q. Yang, Further results for Perron-Frobenius theorem for non-negative tensors, SIAM J. Matrix Anal. Appl. 31(2010), 2517-2530.
* [8] Y Yang and Q. Yang, Further results for Perron-Frobenius theorem for non-negative tensors II, SIAM J. Matrix Anal. Appl. 31(2011), 1236-1250.
|
arxiv-papers
| 2013-09-20T07:11:32 |
2024-09-04T02:49:51.199986
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Pingzhi Yuan, Lihua You",
"submitter": "Lihua You",
"url": "https://arxiv.org/abs/1309.5189"
}
|
1309.5231
|
11institutetext: Institute for Astronomy and Astrophysics, Section
Computational Physics, Eberhard Karls Universität Tübingen, Auf der
Morgenstelle 10, D-72076 Tübingen, Germany
11email: [email protected], [email protected],
[email protected]
22institutetext: Dipartimento di Fisica, Università degli Studi di Torino, via
Pietro Giuria 1, 10125 Torino, Italy
# Radiation hydrodynamics integrated in the code PLUTO
Stefan M. Kolb 11 Matthias Stute 11 Wilhelm Kley 11 Andrea Mignone 22
(Received ; accepted 9. September 2013)
###### Abstract
Aims. The transport of energy through radiation is very important in many
astrophysical phenomena. In dynamical problems the time-dependent equations of
radiation hydrodynamics have to be solved. We present a newly developed
radiation-hydrodynamics module specifically designed for the versatile MHD
code PLUTO .
Methods. The solver is based on the flux-limited diffusion approximation in
the two-temperature approach. All equations are solved in the co-moving frame
in the frequency independent (grey) approximation. The hydrodynamics is solved
by the different Godunov schemes implemented in PLUTO , and for the radiation
transport we use a fully implicit scheme. The resulting system of linear
equations is solved either using the successive over-relaxation (SOR) method
(for testing purposes), or matrix solvers that are available in the PETSc
library. We state in detail the methodology and describe several test cases in
order to verify the correctness of our implementation. The solver works in
standard coordinate systems, such as Cartesian, cylindrical and spherical, and
also for non-equidistant grids.
Results. We have presented a new radiation-hydrodynamics solver coupled to the
MHD-code PLUTO that is a modern, versatile and efficient new module for
treating complex radiation hydrodynamical problems in astrophysics. As test
cases, either purely radiative situations, or full radiation-hydrodynamical
setups (including radiative shocks and convection in accretion discs) have
been studied successfully. The new module scales very well on parallel
computers using MPI. For problems in star or planet formation, we have added
the possibility of irradiation by a central source.
###### Key Words.:
radiation transport – irradiation – hydrodynamics – accretion disc
## 1 Introduction
Radiative effects play a very important role in nearly all astrophysical fluid
flows, ranging from planet and star formation to the largest structures in the
universe. Coupling the equations of radiation transport to those of
(magneto-)hydrodynamics (MHD) has been studied for decades, and comprehensive
treatments can be found for example in text books by Mihalas & Mihalas (1984)
or Pomraning (1973). The numerical implementation of two-temperature radiation
hydrodynamics (in the diffusion approximation) into multi-dimensional MHD/HD-
codes has been done already over twenty years ago in various implementations,
for example by Eggum et al. (1988), Kley (1989), in the ZEUS-code (Stone et
al., 1992), and later by Turner & Stone (2001).
In order to study, for example, the dynamics and characteristics of stellar
atmospheres together with convection, more accurate solvers for the radiation
transport based on the method of short characteristics have been developed,
see Davis et al. (2012) and Freytag et al. (2012) for the present status. This
can then be coupled to the hydrodynamics using the Variable Eddington Tensor
method (Jiang et al., 2012). Another approach is the $\mathrm{M}1$ closure
model where the radiative moment equations are closed at a higher level
(González et al., 2007; Aubert & Teyssier, 2008). Despite this progress it is
still useful and desirable to have a method at hand which solves the
interaction of matter and radiation primarily within the bulk part of the
matter which may be optically thick. In such type of applications, the method
of flux-limited diffusion (FLD, see Levermore & Pomraning (1981)) has its
clear merits and is still implemented into existing MHD-codes, for example in
NIRVANA (Kley et al., 2009) to study the planet formation process, in RAMSES
(Commerçon et al., 2011) for protostellar collapse simulations, and in
combination with a multi-frequency irradiation tool into PLUTO (Kuiper et al.,
2010) for massive star formation.
Since the 3D-MHD code PLUTO (Mignone et al., 2007) is becoming increasingly
popular within the computational astrophysics community, we added a publicly
available radiation module, which is based on the two-temperature FLD-
approximation, as described by Commerçon et al. (2011). PLUTO solves the
equations of hydrodynamics and magnetohydrodynamics including the non-ideal
effects of viscosity, thermal conduction and resistivity by means of shock-
capturing Godunov-type methods. Several Riemann solvers, several time stepping
methods and interpolation schemes can be chosen. Additionally, we added a ray-
tracing routine that allows for additional irradiation by a point source in
the center. Treating the irradiation in a ray-tracing approach, guarantees the
long-range character of the radiation better than FLD (Kuiper et al., 2012;
Kuiper & Klessen, 2013).
The paper is organized as follows. In section 2.1, we briefly introduce the
equations of hydrodynamics including radiation transport. Additionally we
describe the general idea behind the flux limited diffusion approximation. In
section 3, we present the discretization of the equations and the solver of
the resulting matrix equation, and present our numerical implementation of
irradiation. In section 4, we present six different test cases to show the
correctness of the implemented equations: four test cases with an analytical
solution (section 4.1 to 4.4) and two others in which our results are compared
with those from other codes (section 4.5 and 4.6). We end with a summary and
conclusions.
## 2 Radiation hydrodynamics
### 2.1 The equations
Even though the PLUTO-environment includes the full MHD-equations and non-
ideal effects such as viscosity, we restrict ourselves here to the Euler
equations of ideal hydrodynamics. Radiation effects are included in the two-
temperature approximation, which implies an additional equation for the
radiation energy. In order to follow the transport of radiation, we apply the
flux-limited diffusion approximation and treat the exchange of energy and
momentum between the gas and the radiation field with additional terms in the
gas momentum and energy equations. The system of equations then read:
$\displaystyle\frac{\partial}{\partial t}\rho+\nabla\cdot(\rho\@vec{v})$
$\displaystyle=0$ (1) $\displaystyle\frac{\partial}{\partial
t}\rho\@vec{v}+\nabla\cdot\left(\rho\@vec{v}\otimes\@vec{v}\right)+\nabla p$
$\displaystyle=+\rho(\@vec{a}_{\mathrm{ext}}+\@vec{a}_{\mathrm{rad}})$ (2)
$\displaystyle\frac{\partial}{\partial
t}e+\nabla\cdot\left[\left(e+p\right)\@vec{v}\right]$
$\displaystyle=+\rho\@vec{v}\cdot(\@vec{a}_{\mathrm{ext}}+\@vec{a}_{\mathrm{rad}})$
(3) $\displaystyle\phantom{=}\hskip 3.0pt-\kappa_{\mathrm{P}}\rho
c(a_{\mathrm{R}}T^{4}-E)$ $\displaystyle\frac{\partial}{\partial
t}E+\nabla\cdot\@vec{F}$ $\displaystyle=\kappa_{\mathrm{P}}\rho
c\left(a_{\mathrm{R}}T^{4}-E\right)$ (4)
The first three equations (1-3) describe the evolution of the gas motion,
where $\rho$ is the gas density, $p$ the thermal pressure, $\@vec{v}$ the
velocity, $e=\rho\,\epsilon+1/2\,\rho\,v^{2}$ the total energy density (i.e.,
the sum of internal and kinetic) of the gas without radiation, and
$\@vec{a}_{\mathrm{ext}}$ an acceleration caused by external forces (e.g.
gravity), not induced by the radiation field (see below). This system of
equations is closed by the ideal gas relation
$p=(\gamma-1)\,\rho\,\epsilon=\rho\frac{k_{\mathrm{B}}\,T}{\mu\,m_{\mathrm{H}}},$
(5)
where $\gamma$ is the ratio of specific heats, $T$ the gas temperature,
$k_{\mathrm{B}}$ the Boltzmann constant, $\mu$ the mean molecular weight, and
$m_{\mathrm{H}}$ the mass of hydrogen. The specific internal energy can be
written as $\epsilon=c_{\mathrm{V}}\,T$, with the specific heat capacity given
by
$c_{\mathrm{V}}=\frac{k_{\mathrm{B}}}{(\gamma-1)\mu m_{\mathrm{H}}}.$ (6)
Here, we assume constant $\gamma$ and $\mu$, which also implies a constant
$c_{\mathrm{V}}$.
The evolution of the radiation energy density $E$ is given by eq. (4), where
$\@vec{F}$ denotes the radiative flux, $\kappa_{\mathrm{P}}$ the Planck mean
opacity, $c$ the speed of light and $a_{\mathrm{R}}$ the radiation constant.
The fluid is influenced by the radiation in two different ways. First, the
radiation may be absorbed or emitted by the fluid leading to variation of its
energy density. This variation is given by the expression
$\kappa_{\mathrm{P}}\rho c\left(a_{\mathrm{R}}T^{4}-E\right)$, see right hand
side of equation (3) and (4). The second effect is that of radiation pressure.
We include this term as an additional acceleration to the momentum equation,
$\@vec{a}_{\mathrm{rad}}=\frac{\kappa_{\mathrm{R}}}{c}\@vec{F}$. The present
implementation does not include the advective transport terms for the
radiation energy and radiative pressure work in eqs. (3) and (4). For the
relatively low temperature protoplanetary disk application that we consider
here these terms are of minor importance. If required, these terms can be
treated in our implementation straightforwardly within PLUTO by adding
additional source terms.
### 2.2 The flux-limited diffusion approximation
The system of equations shown cannot be solved without further assumptions for
the radiative flux $\@vec{F}$. Here we use the flux limited diffusion
approximation (FLD) where the radiation flux is given by a diffusion
approximation
$\@vec{F}=-\lambda\,\frac{c}{\kappa_{\mathrm{R}}\,\rho}\,\nabla E\,,$ (7)
with the Rosseland mean opacity $\kappa_{\rm R}$. The flux-limiter $\lambda$
describes approximately the transition from very optically thick regions with
$\lambda=1/3$ to optically thin regimes, where $\@vec{F}\rightarrow-
cE\frac{\nabla E}{|\nabla E|}$. This leads to the formal definition of the
flux-limiter which is a function of the dimensionless quantity
$R=\frac{|\nabla E|}{\kappa_{\mathrm{R}}\rho E}\,,$ (8)
with the following behaviour:
$\lambda(R)=\left\\{\begin{array}[]{ll}\frac{1}{3},&R\to 0\\\
\frac{1}{R},&R\to\infty\end{array}\right.$ (9)
Physically sensible flux-limiters thus have to fulfil the equation (9) in the
given limits and describe the behaviour between the limits approximately. We
have implemented three different flux-limiters:
$\displaystyle\lambda(R)$ $\displaystyle=\frac{1}{R}\left(\coth
R-\frac{1}{R}\right)$ (10) $\displaystyle\lambda(R)$
$\displaystyle=\left\\{\begin{array}[]{l l}\frac{2}{3+\sqrt{9+12R^{2}}}&\quad
0\leq R\leq\frac{3}{2}\\\
\frac{1}{1+R+\sqrt{1+2R}}&\quad\frac{3}{2}<R\leq\infty\\\ \end{array}\right.$
(13) $\displaystyle\lambda(R)$ $\displaystyle=\left\\{\begin{array}[]{l
l}\frac{2}{3+\sqrt{9+10R^{2}}}&\quad 0\leq R\leq 2\\\
\frac{10}{10R+9+\sqrt{180R+81}}&\quad 2<R\leq\infty\\\ \end{array}\right.$
(16)
from Levermore & Pomraning (1981), Minerbo (1978), and Kley (1989),
respectively. A comparison of them is presented in Kley (1989).
In general it is necessary to solve the equations for each frequency which
appears in the physical problem. However, here we use the grey approximation
in which all radiative quantities including the opacities are integrated over
all frequencies. In our treatment scattering is not accounted for directly,
but it is included in the effective isotropic absorption and emission
coefficients.
## 3 Solving the radiation part
### 3.1 Reformulation of the equations
Instead of solving system of equations (1-4) directly as a whole, the problem
is split into two steps. In the first step, PLUTO is used to solve the
equations of fluid dynamics with the additional force caused by the radiation.
This corresponds to the equations (1) to (3) with the additional acceleration,
$\@vec{a}_{\mathrm{rad}}$, but without the interaction term between the matter
and radiation (last term in eq. 3). By using PLUTO for solving the non-
radiative part of the equations, we are not limited to the Euler equations,
but are able to use the full capabilities of PLUTO for solving the equations
of hydrodynamics or magnetohydrodynamics, including the effects of viscosity
and magnetic resistivity.
In a second, additional step we solve the radiation energy equation (4) and
for the corresponding heating-cooling term in the internal energy of the
fluid:
$\left.\begin{array}[]{lll}\displaystyle{\frac{\partial}{\partial
t}E-\nabla\cdot\left(\frac{c\lambda}{\kappa_{\mathrm{R}}\rho}\nabla
E\right)}&=&\phantom{-}\kappa_{\mathrm{P}}\rho
c\left(a_{\mathrm{R}}T^{4}-E\right)\\\ \displaystyle{\frac{\partial}{\partial
t}\rho\epsilon}&=&-\kappa_{\mathrm{P}}\rho
c\left(a_{\mathrm{R}}T^{4}-E\right)\end{array}\quad\right\\}$ (17)
In order to obtain the radiation energy density, we solve the system of
coupled equations (17). Within one time step PLUTO advances the hydrodynamical
quantities, i.e. the density $\rho$, the velocity $\@vec{v}$ and a temporary
pressure $p$ from time $t^{n}$ to the time $t^{n+1}$, where the time step,
$\Delta t=t^{n+1}-t^{n}$, is determined by PLUTO using the CFL conditions,
presently without including the radiation pressure. These depend on the used
time stepping method in PLUTO , for more information see Mignone et al. (2007)
and the userguide of PLUTO .
The physical process of radiation transport takes place on time scales much
shorter than the one in hydrodynamics. In order to use the same time step for
hydrodynamics and the radiation transport, we apply an implicit scheme to
handle the radiation diffusion and the coupling between matter and radiation
described by equation (17). Because of the coupling of the equations, the
method will update $T$ and $E$ simultaneously, which leads formally to a
nonlinear set of coupled equations. As outlined below, the system is solved
for the radiation energy density $E$. From the new values for $E$, we compute
the new fluid temperature (see eq. 21 below) and update the fluid pressure by
using the ideal gas relation from equation (5). This is then used within PLUTO
to calculate a new total gas energy $e$.
### 3.2 Discretization
In order to discretize the equations (17), we apply a finite volume method.
For that purpose we integrate over the volume of a grid cell and transform the
divergence into a surface integral. Furthermore, we replace the gradient of
$E$ by finite differences, and apply an implicit scheme. The discretization
scheme has been implemented in 3D for Cartesian, cylindrical and spherical
polar coordinates including all the necessary geometry terms for the
divergence and gradient. Since the density has been updated already in the
hydrodynamical part of the solver, we can replace
$\frac{\partial\rho\epsilon}{\partial t}$ with
$\rho\,c_{\mathrm{V}}\,\frac{\partial T}{\partial t}$, which is valid for a
constant heat capacity. Then the resulting discretized equations for the
radiative part can be written as
$\displaystyle\frac{E_{i,j,k}^{n+1}-E_{i,j,k}^{n}}{\Delta t}$ (18)
$\displaystyle=$ $\displaystyle
G_{x1}^{r}K_{i+\frac{1}{2},j,k}^{n}\frac{E_{i+1,j,k}^{n+1}-E_{i,j,k}^{n+1}}{\Delta{x_{1}}_{i+\frac{1}{2}}}-G_{x1}^{l}K_{i-\frac{1}{2},j,k}^{n}\frac{E_{i,j,k}^{n+1}-E_{i-1,j,k}^{n+1}}{\Delta{x_{1}}_{i-\frac{1}{2}}}$
$\displaystyle+$ $\displaystyle
G_{x2}^{r}K_{i,j+\frac{1}{2},k}^{n}\frac{E_{i,j+1,k}^{n+1}-E_{i,j,k}^{n+1}}{\Delta{x_{2}}_{j+\frac{1}{2}}}-G_{x2}^{l}K_{i,j-\frac{1}{2},k}^{n}\frac{E_{i,j,k}^{n+1}-E_{i,j-1,k}^{n+1}}{\Delta{x_{2}}_{j-\frac{1}{2}}}$
$\displaystyle+$ $\displaystyle
G_{x3}^{r}K_{i,j,k+\frac{1}{2}}^{n}\frac{E_{i,j,k+1}^{n+1}-E_{i,j,k}^{n+1}}{\Delta{x_{3}}_{k+\frac{1}{2}}}-G_{x3}^{l}K_{i,j,k-\frac{1}{2}}^{n}\frac{E_{i,j,k}^{n+1}-E_{i,j,k-1}^{n+1}}{\Delta{x_{3}}_{k-\frac{1}{2}}}$
$\displaystyle+$
$\displaystyle{\kappa_{\mathrm{P}}^{n}}_{i,j,k}\,\rho_{i,j,k}^{n}c\left(a_{\mathrm{R}}(T_{i,j,k}^{n+1})^{4}-E_{i,j,k}^{n+1}\right)\,,$
and for the thermal energy (or temperature, respectively)
$\frac{T_{i,j,k}^{n+1}-T_{i,j,k}^{n}}{\Delta
t}=-\frac{{\kappa_{\mathrm{P}}}_{i,j,k}^{n}\,c}{c_{\mathrm{V}}}\,\left(a_{\mathrm{R}}\left(T_{i,j,k}^{n+1}\right)^{4}-E_{i,j,k}^{n+1}\right)\,.$
(19)
Here, the superscript $n$ refers to the values of all variables after the most
recent update from the hydrodynamical step. In order to simplify the notation
for the separate radiation module, we assume the update takes place from time
$n$ to $n+1$. The subscripts $i,j,k$ refer to the 3 spatial directions of the
computational grid, where all variables are located at the cell centers. Half-
integer indices refer to cell interfaces. The physical sizes (proper length)
of each cell in the 3 spatial directions $m$ ($m=1,2,3$) are given by
$\Delta{x_{m}}$, where we additionally allow for non-equidistant grids. The
effective radiative diffusion coefficient (defined at cell centers) is given
by
$K_{i,j,k}^{n}=\frac{c\lambda(R_{i,j,k})}{{\kappa_{\mathrm{R}}}_{i,j,k}^{n}\,\rho_{i,j,k}^{n}}\,,$
where $R_{i,j,k}$ is calculated from eq. (8) by central differencing. Values
at cell interfaces are obtained by linear interpolation. The factors
$G^{l,r}_{xm}$ are geometrical terms defined, respectively, as the left and
right surface areas divided by the cell volume in the direction given by
$m=1,2,3$. In the recent work by Bitsch et al. (2013b) the difference
equations have been written out in more detail for Cartesian, equidistant
grids. The required opacities are evaluated using the values of $\rho$ and $T$
after the hydrodynamical update at time $t^{n}$.
As mentioned before, equations (17) constitute a set of coupled nonlinear
equations. The non-linear term $(T_{i,j,k}^{n+1})^{4}$ that appears in
equation (19) is linearised using the method outlined in Commerçon et al.
(2011)
$(T_{i,j,k}^{n+1})^{4}=(T_{i,j,k}^{n})^{4}\left(1+\frac{T_{i,j,k}^{n+1}-T_{i,j,k}^{n}}{T_{i,j,k}^{n}}\right)^{4}\approx
4(T_{i,j,k}^{n})^{3}T_{i,j,k}^{n+1}-3(T_{i,j,k}^{n})^{4}\,.$ (20)
Using this approximation, we obtain an equation for computing the new
temperature in terms of the new radiation energy density, $E_{i,j,k}^{n+1}$,
and the old temperatures, $T_{i,j,k}^{n}$
$T_{i,j,k}^{n+1}=\frac{{\kappa_{\mathrm{P}}}_{i,j,k}^{n}c\left(3a_{\mathrm{R}}(T_{i,j,k}^{n})^{4}+E_{i,j,k}^{n+1}\right)\Delta
t+c_{\mathrm{V}}T_{i,j,k}^{n}}{c_{\mathrm{V}}+4{\kappa_{\mathrm{P}}}_{i,j,k}^{n}ca_{\mathrm{R}}(T_{i,j,k}^{n})^{3}\Delta
t}\,.$ (21)
The expression can be substituted into eq. (18) to obtain a linear system of
equations for the new radiation energies $E_{i,j,k}^{n+1}$, that can be solved
using standard matrix solvers, see section 3.4. The new temperature can then
be calculated from eq. (21). We implemented several boundary conditions for
the radiation energy density including periodic, symmetric and fixed value.
### 3.3 Irradiation
In order to couple possible irradiation to the radiation transport equations,
a new source term, $S$, has to be added to the right hand side of the thermal
energy equation in system (17)
$\frac{\partial\rho\epsilon}{\partial t}=-\kappa_{\mathrm{P}}\rho
c(a_{\mathrm{R}}T^{4}-E)+S\,.$ (22)
This results in an additional term, $S_{i,j,k}/(\rho_{i,j,k}c_{\rm V})$, in
Eq. 19, correspondingly in equation (21), and in a modification of the right
hand side of the resulting matrix equation for $E_{i,j,k}^{n+1}$.
For the present implementation, we assume that the irradiating source is
located at the centre of a spherical coordinate system. Therefore it is
straightforward to compute the optical depth $\tau_{i,j,k}$ even for
simulations using parallel computers. Assuming that a ray of light travels
along the radial direction from the origin to the grid cell $i,j,k$ under
consideration, the optical depth from the inner radius $r_{0}$ to the $i$th
grid cell with radius $r_{i}$ can be simply expressed as the integral along
the radial coordinate,
$\tau_{i,j,k}=\int_{r_{0}}^{r_{i}}\kappa_{\star}\rho(r)\,dr\approx\sum_{n=0}^{i}\kappa_{\star}{{}_{n,j,k}}\,\rho_{n,j,k}\Delta
r_{n}$ (23)
where $\Delta r_{n}$ is the radial length of the $n$th grid cell, and
$\kappa_{\star}$ the opacity used for irradiation. For the sake of
readability, we write $\tau_{i}$ instead of $\tau_{i,j,k}$ in the following.
We use $\kappa_{\star}=\kappa_{\rm P}$ in the test case with irradiation
presented in section 4.3. Additionally $\kappa_{\star}$ can be defined by the
user as well as the other opacities. Re-emission of the photons which were
absorbed in the cell volume is handled in our treatment by the heating-cooling
term see equation (17).
The luminosity of the source is given by
$L_{\star}=4\pi R_{\star}^{2}\sigma T_{\star}^{4}\,,$ (24)
where $\sigma$ denotes the Stefan-Boltzmann constant, $T_{\star}$ is the
temperature of the star and $R_{\star}$ its radius. In order to compute the
amount of irradiated energy which is absorbed by a specific grid cell we have
to know the surface area $A$ of a grid cell oriented perpendicular to the
radiation from the star and the flux $f$ at the radius $r$. This surface area
$A$ is given by the expression
$A_{i,j,k}=\int\limits_{\theta_{j}}^{\theta_{j+1}}\int\limits_{\phi_{k}}^{\phi_{k+1}}\,dA=r_{i}^{2}(\phi_{k+1}-\phi_{k})(\cos\theta_{j}-\cos\theta_{j+1})\,,$
(25)
where $\theta$ is the azimuthal and $\phi$ the polar angle in the spherical
coordinate system. Without absorption the flux $f$ is given by the expression
$f=\frac{L_{\star}}{4\pi r^{2}}=\sigma
T_{\star}^{4}\left(\frac{R_{\star}}{r}\right)^{2}\,.$ (26)
The amount of energy per time which arrives at the surface of the grid cell
$(i,j,k)$ is
$H_{i,j,k}=A_{i,j,k}f=(\phi_{k+1}-\phi_{k})(\cos\theta_{j}-\cos\theta_{j+1})\sigma
T_{\star}^{4}R_{\star}^{2}\,,$ (27)
again without absorption. If the irradiated energy is partly absorbed, the
remaining amount of energy per time is then $H_{i,j,k}e^{-\tau_{i,j,k}}$.
Using these results, we can compute the energy density per time, $S$, which is
absorbed by one grid cell ($i,j,k$)
$\displaystyle S_{i,j,k}$
$\displaystyle=\frac{H_{i,j,k}e^{-\tau_{i}}-H_{i,j,k}e^{-\tau_{i+1}}}{V_{i,j,k}}=\frac{H_{i,j,k}\left(e^{-\tau_{i}}-e^{-\tau_{i+1}}\right)}{V_{i,j,k}}$
$\displaystyle=\frac{3\sigma
T_{\star}^{4}R^{2}_{\star}\left(e^{-\tau_{i}}-e^{-\tau_{i+1}}\right)}{(r_{i+1}^{3}-r_{i}^{3})}\,,$
(28)
with the volume of a grid cell
$\displaystyle V_{i,j,k}$
$\displaystyle=\int\limits_{r_{i}}^{r_{i+1}}\int\limits_{\theta_{j}}^{\theta_{j+1}}\int\limits_{\phi_{k}}^{\phi_{k+1}}r^{2}sin\theta\,dr\,d\theta\,d\phi$
$\displaystyle=\frac{1}{3}(r_{i+1}^{3}-r_{i}^{3})(\cos\theta_{j}-\cos\theta_{j+1})(\phi_{k+1}-\phi_{k})\,.$
(29)
The absorbed energy density per time, $S_{i,j,k}$, is computed for each grid
cell before solving the matrix equation. A similar treatment of irradiation
has been described recently by Bitsch et al. (2013b), for a multi-frequency
implementation see Kuiper et al. (2010).
### 3.4 The matrix solver
We implemented two different solvers for the matrix equation. The first one
uses the method of successive over-relaxation (SOR), and as a faster and more
flexible solver we use the PETSc111For more information visit the website
http://www.mcs.anl.gov/petsc or have a look at Balay et al. (2012). library.
From the PETSc library we use the Krylov subspace iterative method and a
preconditioner to solve the matrix equation. For all test cases described we
used gmres (Generalized Minimal Residual) as iterative methode and bjacobi
(Block Jacobi) as preconditioner. Beside others the convergence of the SOR
algorithm and the PETSc library can be estimated using the following criteria
$\left\|\@vec{r}^{(k)}\right\|<\max(\epsilon_{\mathrm{r}}\cdot\left\|\@vec{b}\right\|,\epsilon_{\mathrm{a}})$
(30)
where $\@vec{b}$ is the right hand side of the matrix equation
$A\@vec{x}=\@vec{b}$, $\@vec{r}^{(k)}=\@vec{b}-A\@vec{x}^{(k)}$ is the
residual vector for the k-th iteration of the solver and $\@vec{x}$ is the
solution vector (here the radiation energy density). As norm we used here the
$L_{2}$ norm. The quantities $\epsilon_{\mathrm{r}}$ and
$\epsilon_{\mathrm{a}}$ are the relative and absolute tolerance, respectively,
and are problem dependent, with a common value of $10^{-50}$ for
$\epsilon_{\mathrm{a}}$. For the test cases in section 4 we use relative
tolerances $\epsilon_{\mathrm{r}}$ between $10^{-5}$ and $10^{-8}$. The
criterion (30) is the default one used by the PETSc library. For more
information about the convergence test in PETSc the reader should refer to
section 4.3.2. of Balay et al. (2012). The solver performance in a parallel
environment is described in section 4.6.4.
## 4 Test cases
In order to verify the implemented method, we simulated several test problems
and compared the results with either corresponding analytical solutions or
calculations done with different numerical codes. Most of the tests correspond
to one-dimensional problems. In order to model those, we have used quasi one-
dimensional domains, with a very long cuboid that has the height $h$, width
$w$ and a length $l$. The length $l$ is much larger than the width or height,
and for simplicity we use $w=h$. We performed some of the tests in all three
implemented coordinate systems (Cartesian, cylindrical and spherical) and in
three different alignments of the cuboid along each coordinate direction. This
is done to check whether the geometry factors are correct. In the case of a
non-Cartesian coordinate system we placed the cuboid at large distances $r$
from the origin such that the domain approximately describes a Cartesian
setup.
We use for all test cases the solver based on the PETSc library with the
default iterative solver gmres and the pre-conditioner bjacobi.
### 4.1 Linear diffusion test
Figure 1: Linear diffusion test at the time $t=4.2\cdot 10^{-12}\,\mathrm{s}$.
The simulated (read dots) and the analytical (black line) solution is plotted.
We also plot the solution with the absolute value of the relative error (blue
dashed line) which belongs to the axis on the right. Figure 2: Time evolution
for the linear diffusion test from time $t=0\,\mathrm{s}$ to $4.2\cdot
10^{-12}\,\mathrm{s}$ at three different positions at $x=0\,\mathrm{cm}$
(black lines), at $x=0.5\,\mathrm{cm}$ (blue lines) and at
$x=1.0\,\mathrm{cm}$ (red lines). The dotted lines for each position belong to
the simulated solution, the solid lines to the analytical solution and the
dashed lines show the relative error which belong to the axis on the right.
The following test is adapted from Commerçon et al. (2011). The initial
profile of the radiation energy density is set to a delta function which is
then evolved in time and compared to the analytical one dimensional solution.
We perform this test in all implemented coordinate systems (Cartesian,
cylindrical and spherical coordinates) as described above, which results in
nine different simulations. The used domain is quasi one-dimensional and the
equations of hydrodynamics are not solved in this test. Only the radiation
diffusion equation
$\frac{\partial}{\partial
t}E=\nabla\cdot\left(\frac{c\lambda}{\kappa_{\mathrm{R}}\rho}\nabla E\right)$
(31)
is solved which we obtain from equations (17) by setting
$\kappa_{\mathrm{P}}=0$. An analytical solution to equation (31) can be
calculated in the one dimensional case with a constant flux-limiter
$\lambda=\frac{1}{3}$ and a constant product of the Rosseland opacity and
density, here we set $\kappa_{\mathrm{R}}\rho=1\,\mathrm{cm^{-1}}$. The
equation to solve is then given by
$\frac{\partial}{\partial t}E(x,t)=\frac{c}{3}\frac{\partial^{2}}{\partial
x^{2}}E(x,t)$ (32)
with solution
$E(x,t)=\frac{\tilde{E}_{0}}{\sqrt{\frac{4}{3}c\pi
t}}e^{-\frac{3x^{2}}{4ct}}\,,$ (33)
where $\tilde{E}_{0}$ is the integral over the initial profile of the energy
density, $E(x,t=0)$. Note that in the quasi one-dimensional case (using a
stretched 3D domain) $\tilde{E}_{0}$ has the units $\mathrm{erg\,cm^{-2}}$.
#### 4.1.1 Setup
The domain is a cuboid with a length of $4\,\mathrm{cm}$ and a width and
height of $0.04\,\mathrm{cm}$. We used here $301\times 3\times 3$ grid cells.
The initial profile of the radiation energy density in the quasi one-
dimensional case is set by
$E_{i}=\begin{cases}1\,\mathrm{\frac{erg}{cm^{3}}},&i=1,2,\dots,N\text{ with
}i\neq\frac{N}{2}\\\\[5.0pt] \hskip 10.0pt\frac{\tilde{E}_{0}}{\Delta
x},&i=\frac{N}{2}\\\ \end{cases}$ (34)
where $\Delta x$ is the length of a grid cell. For numerical reasons, we have
set $E_{i}$ for $i\neq\frac{N}{2}$ to the value $1\,\mathrm{erg\,cm^{-3}}$
instead of $0\,\mathrm{erg\,cm^{-3}}$. This choice is not problematic, since
$\tilde{E}_{0}/\Delta x\gg 1\,\mathrm{erg\,cm^{-3}}$ for our chosen value of
$\tilde{E}_{0}=10^{5}\,\mathrm{erg\,cm^{-2}}$. The initial values for pressure
and density are $p=1\,\mathrm{g\,cm^{-1}\,s^{-2}}$ and
$\rho=1\,\mathrm{g\,cm^{-3}}$. Furthermore we use
$\kappa_{\mathrm{R}}=1\,\mathrm{cm^{2}\,g^{-1}}$ for the Rosseland opacity.
All boundary conditions are set to periodic except for the boundary conditions
at the beginning and end of the quasi one-dimensional domain, which are set to
outflow. For the matrix solver we used a relative tolerance of
$\epsilon_{r}=10^{-8}$. The simulation starts at $t=0\,\mathrm{s}$ with an
constant time step of $\Delta t=1\cdot 10^{-14}\,\mathrm{s}$ and stops at
$t=4.2\cdot 10^{-12}\,\mathrm{s}$.
#### 4.1.2 Results
The numerical solution $E_{\mathrm{n}}$ and the analytical solution
$E_{\mathrm{a}}$ from equation (33) are plotted in the figures 1 and 2
together with the absolute value of the relative error. In figure 1 the
radiation energy density is plotted against the position at the time
$t=4.2\cdot 10^{-12}\,\mathrm{s}$. The relative error in the relevant range
from $-1\,\mathrm{cm}$ to $1\,\mathrm{cm}$ is always below one percent. In
figure 2 the time evolution from $t=0\,\mathrm{s}$ to $4.2\cdot
10^{-12}\,\mathrm{s}$ is shown for the positions
$x=\\{0,0.5,1.0\\}\,\mathrm{cm}$ coded in the colors black, blue and red,
respectively. The results shown in this figure depend strongly on the
position. For the position $x=0\,\mathrm{c}m$ the error is, for all times
later than $t=4\cdot 10^{-13}\,\mathrm{s}$, below one percent and decreases
with time. For the other positions, the behaviour is different. The relative
error rises and after a while it decreases. This behaviour can be explained by
looking at figure 1. The error is higher at the diffusion front. This region
moves with time and causes the effect for the other positions. The test shows
that the time evolution of the radiation energy density is reproduced
correctly. As described, this test was performed in different coordinate
systems and orientations, with the same results.
### 4.2 Coupling test
Figure 3: Coupling test from $t=10^{-20}\,\mathrm{s}$ to
$t=10^{-4}\,\mathrm{s}$ with three different initial gas energy densities. The
reference solution (black lines) and the simulated results for the initial
energy density $e_{0}=10^{10}\,\mathrm{erg\,cm^{-3}}$ (red dots),
$e_{0}=10^{6}\,\mathrm{erg\,cm^{-3}}$ (blue dots) and
$e_{0}=10^{2}\,\mathrm{erg\,cm^{-3}}$ (green dots) are plotted.
The purpose of this test from Turner & Stone (2001) is to check the coupling
between radiation and the fluid. For this purpose we simulate a stationary
fluid which is initially out of thermal equilibrium. In this simulation the
radiation energy density is the dominant energy which is constant over the
whole simulation. The system of equations (17) decouples in this case and, in
addition, it is not necessary to solve the matrix equation for $E$. By setting
$\sigma_{\mathrm{P}}=\kappa_{\mathrm{P}}\rho$ and $T=\frac{p}{\rho}\frac{\mu
m_{\mathrm{H}}}{k_{\mathrm{B}}}$ from eq. (5) with the assumption that
$\sigma_{\mathrm{P}}$ and $\rho$ are constant, we can rewrite the thermal
energy equation of the system (17) as
$\frac{\,\mathrm{d}e}{\,\mathrm{d}t}=\underbrace{c\sigma_{\mathrm{P}}E}_{C_{\mathrm{1}}}-\underbrace{c\sigma_{\mathrm{P}}a_{\mathrm{R}}\left(\frac{\gamma-1}{\rho}\frac{\mu
m_{\mathrm{H}}}{k_{\mathrm{B}}}\right)^{4}}_{C_{\mathrm{2}}}e^{4}\,.$ (35)
With the used approximations, the coefficients $C_{\mathrm{1}}$ and
$C_{\mathrm{2}}$ are constant. The solution to Eq. (35) can be calculated
analytically in terms of an algebraic equation which would have to be solved
iteratively. Hence, we integrate Eq. (35) numerically using a Runge Kutta
solver of $4$-th order scheme with adaptive step size. In the following we
refer to this solution as the reference solution.
Information on the expected behaviour of the solution can be obtained directly
from the differential equation. It is clear that in the final equilibrium
state (with $\frac{\,\mathrm{d}e}{\,\mathrm{d}t}=0$) the gas temperature has
to be equal to the radiation temperature
$T=\sqrt[4]{\frac{E}{a_{\mathrm{R}}}}$, thus the final gas energy density will
be
$e_{\mathrm{final}}=\left(\frac{C_{\mathrm{1}}}{C_{\mathrm{2}}}\right)^{\frac{1}{4}}\,.$
(36)
If the initial gas energy density $e_{0}$ is much lower than
$e_{\mathrm{final}}$, we can neglect the second term in eq. (35) at the
beginning, thus $e(t)=C_{\mathrm{1}}\,t+e_{0}$. The corresponding coupling
time can be estimated to
$\tau=\frac{e_{\mathrm{final}}-e_{0}}{C_{\mathrm{1}}}\,.$ (37)
On the other hand, if $e_{0}\gg e_{\mathrm{final}}$, we can neglect the first
term in eq. (35) and derive
$e(t)\propto\left(C_{\mathrm{2}}\,t\right)^{-\frac{1}{3}}\quad\mathrm{and}\quad\tau=\frac{1}{e_{\mathrm{final}}^{3}\,C_{\mathrm{2}}}\,.$
(38)
#### 4.2.1 Setup
The computational domain is identical to that of the linear diffusion test in
section 4.1. For the grid we use a resolution of $25\times 3\times 3$ grid
cells. As before we do not solve the equations of hydrodynamics and the
boundary conditions are quite simple. All boundaries are set to periodic
boundary conditions. The constants we used are set to: radiation energy
density $E=10^{12}\,\mathrm{erg\,cm^{-3}}$, density
$\rho=10^{-7}\,\mathrm{g\,cm^{-3}}$, opacity $\sigma_{\mathrm{P}}=4\cdot
10^{-8}\,\mathrm{cm^{-1}}$, mean molecular weight $\mu=0.6$ and the ratio of
specific heats $\gamma=5/3$. The simulations starts at $t=0\,\mathrm{s}$ with
an initial time step of $\Delta t=10^{-20}\,\mathrm{s}$ and evolves until
$t=10^{-4}\,\mathrm{s}$. After each step the time step is increased by 1% in
order to speed-up the computation. The simulation is done with three different
initial gas energy densities, $e_{0}=10^{10}\,\mathrm{erg\,cm^{-3}}$,
$e_{0}=10^{6}\,\mathrm{erg\,cm^{-3}}$ and
$e_{0}=10^{2}\,\mathrm{erg\,cm^{-3}}$.
#### 4.2.2 Results
Figure 3 shows the numerical gas energy density and the reference solution
plotted against time for the three different initial values of $e$. The
agreement of both results is excellent for all initial values. From the figure
we see that in the limit of small and large initial $e_{0}$, we find exactly
the behaviour as predicted by the estimates for eq. (35). The analytic
estimates for the coupling time $\tau$ from equation (37) agrees very well
with our results. The estimate for $e_{0}=10^{2}\,\mathrm{erg\,cm^{-3}}$ is
$\tau=5.88\cdot 10^{-8}\,\mathrm{s}$ and for
$e_{0}=10^{6}\,\mathrm{erg\,cm^{-3}}$ we calculated $\tau=5.78\cdot
10^{-8}\,\mathrm{s}$. In the case that $e_{\mathrm{final}}<e_{0}$ the estimate
in eq. (38) is approximate $\tau=5.88\cdot 10^{-8}\,\mathrm{s}$. We have to
mention here that this test verifies primarily the correctness of equation
(21). As in the linear diffusion test, this test was performed in three
different coordinate systems in different orientations, with the same results.
### 4.3 Coupling test with irradiation
Figure 4: Coupling test with enabled irradiation from $t=10^{-7}\,\mathrm{s}$
to $t=10^{3}\,\mathrm{s}$ at three different distances $d$ from the inner
boundary of the domain. The reference solution (black lines) and the simulated
results for the energy density $e=10^{2}\,\mathrm{erg\,cm^{-3}}$ are plotted
at distance $d=0\,\mathrm{cm}$ (red dots), $d=3\cdot 10^{4}\,\mathrm{cm}$
(blue dots) and $d=3\cdot 10^{5}\,\mathrm{cm}$ (green dots). Figure 5: Radial
dependency of the gas energy density for the coupling test with enabled
irradiation. The gas energy density is plotted at different times,
$t=1.12\cdot 10^{-4}\,\mathrm{s}$ (red dots), $t=1.62\cdot
10^{-2}\,\mathrm{s}$ (blue dots), $t=1.43\,\mathrm{s}$ (green dots),
$t=1.72\cdot 10^{1}\,\mathrm{s}$ (yellow dots) and $t=9.18\cdot
10^{2}\,\mathrm{s}$ (magenta dots). The dots represent the numerical solution
and the solid black lines the reference solution. The position $d$ is again
measured relative to the inner boundary of the quasi one-dimensional domain.
This test is in its basic setup the same as that from section 4.2, but with
irradiation enabled, i.e. equation (22) is solved instead of the second
equation in (17). As described in section 3.3, irradiation is limited to
spherical coordinates, which we use for this test. With the same assumptions
as in section 4.2 , i.e., that $\sigma_{\mathrm{P}}$ and $\rho$ are constant
and with the definitions for $\sigma_{\mathrm{P}}$, $e$, $p$ as well as for
$T$, it is possible to rewrite $S$ from equation (28) to
$S(r)=\frac{3\sigma
T_{\star}^{4}R^{2}_{\star}e^{-\sigma_{\mathrm{P}}\left(r-r_{\mathrm{0}}\right)}\left(1-e^{-\sigma_{\mathrm{P}}\Delta
r}\right)}{\left(r+\Delta r\right)^{3}-r^{3}}\,,$ (39)
and obtain for equation (22)
$\frac{\,\mathrm{d}e}{\,\mathrm{d}t}=\underbrace{S(r)+c\sigma_{\mathrm{P}}E}_{C_{\mathrm{1}}(r)}-\underbrace{c\sigma_{\mathrm{P}}a_{\mathrm{R}}\left(\frac{\gamma-1}{\rho}\frac{\mu
m_{\mathrm{H}}}{k_{\mathrm{B}}}\right)^{4}}_{C_{\mathrm{2}}}e^{4}\,.$ (40)
The reference solution is computed in the same way as before although it now
depends on the distance $r$ from the star. The quasi one-dimensional domain
starts at $r=9000\cdot 10^{5}\,\mathrm{cm}$ and ends at $r=9003\cdot
10^{5}\,\mathrm{cm}$ and we use $300\times 3\times 3$ grid cells. The domain
size in $\theta$ and $\phi$ direction was chosen in a way such that the grid
cells are nearly quadratic. For the simulation we use a constant radiation
energy density of $E=10^{-2}\,\mathrm{erg\,cm^{-3}}$, a density of
$\rho=10^{-5}\,\mathrm{g\,cm^{-3}}$ , a Rosseland opacity of $\kappa_{\rm
R}=10\,\mathrm{cm^{2}\,g^{-1}}$ and a Planck opacity of
$\kappa_{\mathrm{P}}=\kappa_{\mathrm{R}}$ which corresponds to
$\sigma_{\mathrm{P}}=10^{-4}\,\mathrm{cm^{-1}}$. The opacity for the
irradiation $\kappa_{\star}$ is set to $\kappa_{\mathrm{P}}$. For the star,
the temperature was set to $T_{\star}=6000\,\mathrm{K}$ and the radius to
$R_{\star}=8.1\cdot 10^{8}\,\mathrm{cm}$. Additionally we make the assumption
that there is no absorption in the region between the surface of the star and
the inner boundary of the computation domain. Figure 4 shows the gas energy
density plotted against time with an initial gas energy of
$e=10^{2}\,\mathrm{erg\,cm^{-3}}$ at three different positions
$d=0\,\mathrm{cm}$, $d=3\cdot 10^{4}\,\mathrm{cm}$ and $d=3\cdot
10^{5}\,\mathrm{cm}$ where $d$ is measured relative to the inner boundary of
the quasi one-dimensional domain. In figure 5 the radial dependency of the gas
energy density is plotted for the same simulation at five different times. As
expected, the results show that the gas energy density at a time later than
$t=10^{2}\,\mathrm{s}$ becomes constant and depends on the distance from the
star. The simulated and reference solution show a excellent agreement.
### 4.4 A steady state test
Figure 6: Comparison between the numerical (red dots) and the analytical
(black line) solution of the steady state test after $t=1200\,\mathrm{s}$. In
addition the absolute value of the relative error (blue dashed line) is
plotted related to the axis on the right side. Note this axis is logarithmic.
The original version of this test was published in Flaig (2011). We consider a
one-dimensional stationary setup with a given density stratification. In the
steady state, the time derivatives in the equations (17) vanish and the system
is reduced to the following equation for the radiation energy density
$0=\nabla\cdot\,\left(\frac{c\lambda}{\kappa_{\mathrm{R}}\rho}\nabla
E\right)\,.$ (41)
A further reduction is obtained when we rewrite this equation in one dimension
along the z-axis in Cartesian coordinates. The equation is then much simpler
and can be written as
$\frac{\,\mathrm{d}}{\,\mathrm{d}z}\left(\frac{c\lambda}{\kappa_{\mathrm{R}}\rho}\frac{\,\mathrm{d}}{\,\mathrm{d}z}E\right)=0\,.$
(42)
In general the expression $\frac{c\lambda}{\kappa_{\mathrm{R}}\rho}$ is not
known analytically for realistic opacities. In order to circumvent this
problem, we define the effective optical depth
$\tau_{\mathrm{eff}}=\int\,\mathrm{d}\tau_{\mathrm{eff}}=\int_{z_{b}}^{z_{a}}\kappa_{\mathrm{eff}}\rho\,\mathrm{d}z$
where $z_{a}$ and $z_{b}$ are the lower and upper boundaries of the quasi one-
dimensional domain, respectively, and $\kappa_{\mathrm{eff}}$ is the effective
opacity given by
$\kappa_{\mathrm{eff}}=\frac{1}{3}\frac{\kappa_{\mathrm{R}}}{\lambda}$. By
using
$\,\mathrm{d}\tau_{\mathrm{eff}}=\kappa_{\mathrm{eff}}\rho\,\mathrm{d}z$,
equation (42) can be rewritten as:
$\frac{\,\mathrm{d}}{\,\mathrm{d}z}\frac{\,\mathrm{d}}{\,\mathrm{d}\tau_{\mathrm{eff}}}E=0\,.$
(43)
The solution of this equation is then given by
$E=\left(E(\tau_{\mathrm{eff}}=1)-E(\tau_{\mathrm{eff}}=0)\right)\tau_{\mathrm{eff}}+E(\tau_{\mathrm{eff}}=0)$
(44)
where $E(\tau_{\mathrm{eff}}=0)$ and $E(\tau_{\mathrm{eff}}=1)$ are the
radiation energy density at the position where the effective optical depth has
the values zero or one, respectively. Thus, in the static case the radiation
energy has a linear dependence on the optical depth $\tau_{\mathrm{eff}}$ for
all opacity laws.
#### 4.4.1 Initial setup
The domain was chosen to have an arbitrary length of $300\,\mathrm{cm}$ and a
width and height of $3\,\mathrm{cm}$ and $300\times 3\times 3$ grid cells were
used. This test is performed without solving the hydrodynamical equations,
instead we solved equations (17) for a fixed density and opacity law, and
evolved the solution, until a stationary state has been reached. For the
radiation boundary conditions, we used boundary conditions with fixed values
of $E$ at the lower and upper boundary of the domain. At the lower boundary we
have chosen $E=a_{\mathrm{R}}T^{4}$ with a temperature of
$T=2000\,\mathrm{K}$.
Because the stratification is optically thin at the upper boundary we want to
allow the radiation to escape freely from the domain. For this reason we
simply set the temperature to a very small value at the upper boundary, here
$T=10\,\mathrm{K}$.
All other boundary conditions have been set to periodic. The density
stratification is given by
$\rho(z)=\rho_{0}e^{\frac{1}{2}\left(\frac{z-z_{a}}{0.46\cdot\left(z_{b}-z_{a}\right)}\right)^{2}}\,.$
(45)
The initial temperature profile can be chosen randomly in principle, but in
order to speed up the computation we used a linear temperature profile
starting at $z_{a}$ with $T=2000\,\mathrm{K}$ and ending at $z_{b}$ with
$T=10\,\mathrm{K}$. From this temperature profile we assigned pressure values
using equation (5). The radiation energy density $E$ inside the domain is also
set using the gas temperature profile and $E=a_{\mathrm{R}}T^{4}$. The ratio
of specific heats and the mean molecular weight are set to $\gamma=1.43$ and
$\mu=0.6$, respectively. As flux-limiter we have chosen equation (13), for the
Rosseland mean opacity $\kappa_{\rm R}$ we use data from Lin & Papaloizou
(1985), and the Planck mean opacity is set to $\kappa_{\rm P}=\kappa_{\rm R}$.
The initial time step is $\Delta t=0.3\,\mathrm{s}$ and it is increased
slightly with time in order to speed-up the computation and to keep the number
of iterations done by the matrix solver nearly constant. This simulation was
preformed with a relative tolerance of $\epsilon_{r}=10^{-6}$ for the matrix
solver.
#### 4.4.2 Results
A steady state is reached approximately after $t=1200\,\mathrm{s}$. In figure
6, we plot the radiation energy density against the effective optical depth
$\tau_{\text{eff}}$ from our numerical solution (red dots) together with the
analytical solution from equation (44). The parameters
$E(\tau_{\mathrm{eff}}=1)-E(\tau_{\mathrm{eff}}=0)$ and
$E(\tau_{\mathrm{eff}}=0)$ have been obtained by fitting equation (44) to the
numerical solution. We have to note here that $E(\tau_{\mathrm{eff}}=0)$ is
determined by interpolation between ghost cells and active cells near the
upper boundary $z_{\mathrm{b}}$. Hence, the radiation temperature in the
active region can be much larger than $10\,\mathrm{K}$, a value was
specifically chosen to be very small. We also plot the absolute value of the
relative error $\left|(E_{n}-E_{a})/E_{a}\right|$. The results from the
simulation agree very well with the analytical prediction. As we can see from
figure 6, the largest deviation from the analytical solution is at small
values of $\tau_{\mathrm{eff}}$ with an relative error around one percent. As
the linear diffusion and the coupling test, this test was performed in all
three coordinate systems and in different orientations, with the same results.
### 4.5 Radiation shock
(a) Subcritical shock
(b) Supercritical shock
Figure 7: Sub- and supercritical shock test. In both cases we plot the
radiation temperature (blue line) and the gas temperature (red line) against
$s=z-v\cdot t$ where $z$ is the position along the quasi one-dimensional
domain and $v$ the piston velocity. The subcritical shock 7(a) is shown at
time $t=3.8\cdot 10^{4}\,\mathrm{s}$ and the supercritical shock 7(b) at
$t=7.5\cdot 10^{3}\,\mathrm{s}$.
In this section we extend the previous tests and solve now the full equations
of hydrodynamics and radiation transport simultaneously, testing the complete
new module within PLUTO environment.
#### 4.5.1 Initial setup
Following a set-up from Ensman (1994), a shock is generated in a quasi one-
dimensional domain. This test case is more complex than the ones before, and
it is not possible to derive an analytical solution. Instead we compare our
results with the simulations of Commerçon et al. (2011). The computational
domain is chosen to have a length of $7\cdot 10^{10}\,\mathrm{cm}$ and a width
and height of $3.418\cdot 10^{7}\,\mathrm{cm}$ with $2048\times 4\times 4$
grid cells. The initial density and temperature are set to $\rho=7.78\cdot
10^{-10}\,\mathrm{g\,cm^{-3}}$ and $T=10\,\mathrm{K}$. The initial radiation
energy density is set by the equation $E=a_{\mathrm{R}}T^{4}$. For the flux-
limiter we employ the Minerbo-formulation according to eq. (13), and for the
opacity we use
$\kappa_{\mathrm{R}}\cdot\rho=\kappa_{\mathrm{P}}\cdot\rho=3.1\cdot
10^{-10}\,\mathrm{cm^{-1}}$. Furthermore the ratio of specific heats is set to
$\gamma=7/5$ and the mean molecular weight to $\mu=1$, in analogy to Commerçon
et al. (2011). The time step is computed through the CFL condition of PLUTO
for wich we assume a value of $0.4$. For the solver we took the not so
accurate but robust tvdlf which uses a simple Lax-Friedrichs scheme. For
generating the radiative shock, the following boundary conditions are used: in
the direction of the shock propagation, we employ a reflective boundary
condition at the lower boundary and a zero-gradient at the upper boundary of
the domain. The remaining boundaries are set to periodic. For the relative
tolerance used by the matrix solver we have chosen a value of
$\epsilon_{r}=10^{-5}$. The shock is generated by applying an initial velocity
$v$ to the gas. The velocity is directed towards the reflecting boundary
condition which acts as a wall. The shock propagates then from the wall back
into the domain. Depending on the velocity, the shock is sub- or
supercritical, i.e., the temperature behind the shock front is larger or equal
than the temperature upstream (in front of the shock front), respectively. In
this test we simulate both cases: the subcritical shock with a velocity of
$v=6\cdot 10^{5}\,\mathrm{cm\,s^{-1}}$ and the supercritical shock with
$v=20\cdot 10^{5}\,\mathrm{cm\,s^{-1}}$.
#### 4.5.2 Results
For a better comparison with the results of simulations, where the material is
at rest and a moving piston causes the shock, we introduce the quantity $s$.
This quantity is given by the relation $s=z-v\cdot t$ where $z$ is the
position along the quasi one-dimensional domain. Note that this quantity is
called $z$ in Commerçon et al. (2011). Fig. 7 shows the radiation temperature
(blue line) and the gas temperature (red line) against the previously defined
quantity $s$ for both the subcritical (at $t=3.8\cdot 10^{4}\,\mathrm{s}$) and
supercritical case (at $t=7.5\cdot 10^{3}\,\mathrm{s}$). In the supercritical
case the pre- and post-shock gas temperature are equal, as expected. In the
subcritical case these temperatures can be estimated analytically (Ensman,
1994; Mihalas & Mihalas, 1984; Commerçon et al., 2011). In table 1, the
analytical estimates and the numerical values from our simulations and the
results from Commerçon et al. (2011) are shown together. Here $T_{2}$ is the
post-shock temperature, $T_{-}$ the pre-shock temperature and $T_{+}$ the
spike temperature. In the equations, $R_{\rm G}=\frac{k_{\mathrm{B}}}{\mu
m_{\mathrm{H}}}$ is the perfect gas constant,
$\sigma_{\mathrm{SB}}=\frac{ca_{\mathrm{R}}}{4}$ the Stefan-Boltzmann
constant, and $u$ is the velocity of the shock relative to the upstream
material (or vice versa) in our case $u=7.19\cdot
10^{5}\,\mathrm{cm\,s^{-1}}$.
| analytical estimate | numerical solution | Commerçon et al.
---|---|---|---
$T_{2}\approx\frac{2(\gamma-1)u^{2}}{R_{\rm G}(\gamma+1)^{2}}$ | $\sim\phantom{1}865\,\mathrm{K}$ | $\phantom{1}816.6\,\mathrm{K}$ | $\phantom{1}825\,\mathrm{K}$
$T_{-}\approx\frac{\gamma-1}{\rho uR_{\rm G}}\frac{2\sigma_{\mathrm{SB}}T_{2}^{4}}{\sqrt{3}}$ | $\sim\phantom{1}315\,\mathrm{K}$ | $\phantom{1}331.9\,\mathrm{K}$ | $\phantom{1}275\,\mathrm{K}$
$T_{+}\approx T_{2}+\frac{3-\gamma}{\gamma+1}T_{-}$ | $\sim 1075\,\mathrm{K}$ | $1147.1\,\mathrm{K}$ | $1038\,\mathrm{K}$
Table 1: Comparison of the results from the radiation shock test with
analytical estimates and the results from Commerçon et al. (2011) for the pre-
shock $T_{-}$ and post-shock $T_{2}$ gas temperature as well as the spike
temperature $T_{+}$.
The results agree in general with the analytical estimates and the results
from Commerçon et al. (2011). The analytical estimate for the post-shock
temperature is higher than the numerical results with both codes. We have to
note here that the analytical estimate depends on $u$ and differs therefore
from the values given in Commerçon et al. (2011). The pre-shock and spike
temperatures agree reasonably well with the analytical estimates in our
simulations but are higher than the results from Commerçon et al. (2011). The
differences of our numerical solution to the analytical estimates might due to
the fact that we ignored the advective terms in the radiation energy density
in eq. (4) that may play a role in this dynamic situation. Additionally, it is
noteworthy that the position of the shock front is very well reproduced. This
test was performed in Cartesian coordinates.
### 4.6 Accretion disc
The goal of this last test is to compare the results of different codes on a
more complex two-dimensional physical problem that involves the onset of
convective motions. For this purpose we model a section of an internally
heated, viscous accretion disc in spherical coordinates $(r,\theta,\phi)$
where $r$ is the distance to the centre of the coordinate system, $\theta$ the
polar angle measured from the $z$-axis in cylindrical coordinates and $\phi$
the azimuth angle. The setup follows the standard disc model used in Kley et
al. (2009). The tests proceed in two steps. In a first setup we reduce the
complexity of the problem and consider a static problem, i.e., without solving
the equations of hydrodynamics. This will demonstrate that the equilibrium
between viscous heating and radiative cooling is treated correctly in our
implementation. In the second setup we consider the full hydrodynamic problem
and study the onset of convection in discs.
#### 4.6.1 The initial setup
For both, the static and the dynamical case we use the same initial setup. The
radial extent ranges from $r_{\text{min}}=0.4$ to $r_{\text{max}}=2.5$, where
all lengths are given in units of the semi-major axis of Jupiter
$a_{\text{jup}}=5.2\,\mathrm{AU}$. In the vertical direction the domain
extends from $\theta_{\text{min}}=83\,\mathrm{{}^{\circ}}$ to
$\theta_{\text{max}}=90\,\mathrm{{}^{\circ}}$ and in $\phi$ direction from
$\phi_{\text{min}}=0\,\mathrm{{}^{\circ}}$ to
$\phi_{\text{max}}=360\,\mathrm{{}^{\circ}}$. In the three coordinate
directions $(r,\theta,\phi)$ we use $256\times 32\times 4$ grid cells. The
disc aspect ratio $h$ is set to $h=\frac{H}{s}=0.05\,\mathrm{}$ where
$s=r\sin\theta$ describes the (radial) distance from the $z$-axis in
cylindrical coordinates, and $H$ is the disc’s vertical scale height. The
viscosity $\nu$ is set to a value of $\nu=10^{15}\,\mathrm{cm^{2}\,s^{-1}}$,
and the mean molecular weight to $\mu=2.3\,\mathrm{}$. For the ratio of
specific heats we have used different values, as specified below. The density
stratification can be obtained from vertical hydrostatic equilibrium, assuming
a temperature that is constant on cylinders, $T=T(s)$. It follows (Masset et
al., 2006)
$\rho(r,\theta)=\rho_{0}\cdot
s^{-1.5}\exp\left(\frac{\sin{\theta}-1}{h^{2}}\right)$ (46)
where the quantity $\rho_{0}$ was chosen such that the total mass of the disc
is $M_{\text{disc}}=0.01\cdot M_{\star}$, where $M_{\star}$ is the mass of the
central star of the system which is set to the mass of the sun,
$M_{\star}=M_{\odot}$. The mass within the computational domain is then
$1/2M_{\text{disc}}$ because we only compute the upper half of the disc. The
radial variation leads to a surface density profile of $\Sigma\propto
r^{-1/2}$, which is the equilibrium profile for constant viscosity, and
vanishing mass flux through the disc. The pressure $p$ is set by the
isothermal relation $p=\rho c_{\mathrm{s}}^{2}$, with the speed of sound
$c_{\mathrm{s}}=H\Omega_{\mathrm{K}}$ and the Keplerian angular velocity
$\Omega_{\mathrm{K}}=\sqrt{\frac{GM_{\star}}{s^{3}}}\,,$
with the gravitational constant $G$. The temperature can be computed through
equation (5) and results in $T=\frac{\mu
m_{\mathrm{H}}}{k_{\mathrm{B}}}\frac{p}{\rho}$. The initial velocities are set
to zero except for the angular velocity $v_{\phi}$ which is set to
$v_{\phi}=\sqrt{\frac{\left(1-2h^{2}\right)GM_{\star}}{s}}\,.$
For the Rosseland mean opacity $\kappa_{\rm R}$ we use data from Lin &
Papaloizou (1985), and the Planck mean opacity is set to $\kappa_{\rm
P}=\kappa_{\mathrm{R}}$. The displayed simulations have been performed in the
rotating frame in which the coordinate system rotates with the constant
angular velocity of $\Omega_{\mathrm{K}}$ at $a_{\text{jup}}$, but for non-
rotating systems identical results are obtained. As before the radiation
energy density is initialised to $E=a_{\mathrm{R}}T^{4}$.
For density, pressure and radial velocity we apply reflective radial boundary
conditions and the angular velocity is set to the Keplerian values. In the
azimuthal direction periodic boundary conditions are used for all variables.
In the vertical direction we apply an equatorial symmetry and reflective
boundary condition for $\theta_{min}$. The radiation boundary conditions are
set to reflective for the $r$ direction (both lower and upper), in
$\theta$-direction we use a fixed value of $E=a_{\mathrm{R}}T^{4}$ with
$T=5\,\mathrm{K}$ at $\theta_{min}$ (which denotes the disc surface), and a
symmetric boundary condition holds at the disc’s midplane $\theta_{max}$. For
the $\phi$-direction we use periodic boundary conditions.
In both cases we used for the matrix solver a relative tolerance of
$\epsilon_{r}=10^{-8}$. In the simulation with hydrodynamics we use the
Riemann-solver hllc222Harten, Lax, Van Leer approximate Riemann Solver with
the contact discontinuity.
(a) $t=10\text{ orbits}$
(b) $t=100\text{ orbits}$
Figure 8: Radial mid-plane temperature profile in the simulations with PLUTO
(red dots) and with the code RH2D (black line) after $t=10\,\mathrm{orbits}$
8(a) and $t=100\,\mathrm{orbits}$ 8(b), together with the absolute value of
the relative error (blue dashed line) which belongs to the log axis on the
right.
#### 4.6.2 The static case
In this test case only the radiative equations are solved without the
hydrodynamics. In order to account for the viscous heating in this case, we
add an additional dissipation contribution, $D$, to the right hand side of the
internal energy equation in (17). We consider standard viscous heating, and
include only the main contribution due to the approximately Keplerian shear
flow. At the individual grid points the dissipation is then given by
$D_{i,j,k}=r_{i}^{2}\rho_{i,j,k}\nu\left(\frac{\partial\Omega_{i,j,k}}{\partial
r_{i}}\right)^{2}\,,$ (47)
where $\nu$ is the constant viscosity and $\Omega_{i,j,k}$ the angular
velocity at the individual grid points. In summary we solve the same equations
as in the case with irradiation, when we substitute $S_{i,j,k}$ with
$D_{i,j,k}$.
(a) $\gamma=5/3$
(b) $\gamma=1.1$
Figure 9: Radial mid-plane temperature profile in the simulation with PLUTO
(red line), RH2D (black line) and NIRVANA (blue line) in the quasi-equilibrium
state after 100 orbits in the case with $\gamma=5/3$ without convection 8(a)
and in the strongly convective case with $\gamma=1.1$ 8(b). Additionally we
added the results of a simulation performed with PLUTO where we use a
logarithmic grid in $r$-direction (green line).
In the steady state, the time derivatives in the equations (17) vanish and the
system is reduced to the following equation for the radiation energy density
$\nabla\cdot\,\left(\frac{c\lambda}{\kappa_{\mathrm{R}}\rho}\nabla
E\right)=D.$ (48)
In optically thick regions, $E=a_{\rm R}T^{4}$ and eq. (48) determines the
temperature stratification within the disc.
Figure 10: Vertical slice of the disc temperature at $t=100\,\mathrm{orbits}$
in the dynamical case with $\gamma=1.1$ showing convection cells. Also plotted
in the inset is the enlarged region from $r=0.4\,a_{\text{jup}}$ to
$0.6\,a_{\text{jup}}$ with the velocity field in the $r-\theta$ plane (black
arrows).
The simulation starts at $t=0\,\mathrm{orbits}$ and is evolved until
$t=100\,\mathrm{orbits}$ are reached, where one orbit corresponds to the
Keplerian orbital period at the distance of $a_{\text{jup}}$ which is given
here by $3.732\cdot 10^{8}\,\mathrm{s}$. The initial and overall time step was
chosen as $\Delta t=10^{-3}\,\mathrm{orbits}=3.732\cdot 10^{5}\,\mathrm{s}$.
The results for the static case are shown in figure 8 using here a value of
$\gamma=7/5$ for the adiabatic index. The plots show the radial temperature
profile of the accretion disc in the mid-plane for the simulations after 10
orbits (top panel) and after 100 orbits (bottom panel). We display results of
two different simulations, one done with the code PLUTO (red dots) using the
described methods, and the second (black lines) run with the code RH2D (Kley,
1989). The result from both codes are nearly identical. Even after
$100\,\mathrm{orbits}$ the absolute value of the relative error is always less
than $2\%$. The test shows that the time-scale of the radiative evolution, as
well as the equilibrium state is captured correctly. We note that the code
RH2D uses the one-temperature approach of radiation transport in this case.
#### 4.6.3 The dynamical case
The final equilibrium of the described static case does not depend on the
magnitude of $\gamma$, because the viscous heating is independent of it, see
eq. (48). The situation is different, however, for the dynamical cases, where
the hydrodynamical evolution of the flow is taken into account. Since the time
scale of the radiative transport depends on $\gamma$ (through eq. 6), one
might expect the possibility of convective instability, see for example the
recent work by Bitsch et al. (2013a). This is indeed the case for small enough
values of $\gamma$. In order to demonstrate the correctness of our
implementation also for the full dynamical problem, we modelled two discs, one
with $\gamma=5/3$ which clearly shows no convection, and the other with
$\gamma=1.1$ which shows strong convection. The initial setup was identical to
that described before, but now we solve the equations of viscous hydrodynamics
with radiation transport, but without irradiation and explicit dissipation.
Please note that for viscous flows the energy generation due to viscous
dissipation is automatically included in the total energy equation. The
equations (1) to (3) are solved by PLUTO , and the system of equations (17)
are solved as described in section 3. Since this setup is very dynamical and
requires a more complex interplay of hydrodynamics and radiative transport, we
use an additional third code, NIRVANA, for comparison. The NIRVANA code has
been used in Kley et al. (2009) and Bitsch et al. (2013a) on very similar
setups. The results of the two cases are shown in Fig. 9. In the top panel (a)
we display the result for the $\gamma=5/3$ case which is not convective. Here,
the agreement between the codes is excellent with the maximum deviation in the
percentage range. In the lower panel (b) we display the results for the
$\gamma=1.1$ case. Here the radiative transport time-scale is enhanced which
leads to a strongly convective situation, which can be seen in the raggedness
of the curves. In this simulation we doubled the spatial resolution, compared
with the $\gamma=5/3$ case, such that the convection cells are reasonably well
resolved, see figure 10. The agreement between the three different codes is
very good, despite of the very different solution methods for the
hydrodynamics equations: PLUTO uses the total energy equation with a Riemann-
solver while RH2D and NIRVANA use a second-order upwind scheme and the thermal
energy equation. Additionally, the latter two codes use the full dissipation
function and the one-temperature approach.
#### 4.6.4 Parallel scaling
Figure 11: Parallel scaling benchmark results for the static accretion disc
test case. We plot here the number of processor cores against
$\frac{t_{64}}{t_{\mathrm{N}}}$ where $t_{N}$ is the runtime used on $N$
processors accordingly for $t_{64}$. The used run-times with full
hydrodynamics and radiation transport for $64$, $128$, $256$, $512$ and $1024$
cpu cores (red crosses) are shown together with the ideal case (black dashed
line).
In order to test the parallel scaling of our new implementation, we used the
same setup as in section 4.6.3 and increased the number of grid cells to
$1024\times 64\times 256$. The computations were only run until
$t=5\,\mathrm{orbits}$, and we used the solver PETSc. So we were able to run
the test on $64$ up to $1024$ processor cores within a reasonable time. The
simulations were run on clusters of the BWGrid which are equipped with Intel
Xeon E5440 cpus and have a low latency InfiniBand network. In figure 11 we
show the results of the simulations performed with full hydrodynamics and
radiation transport. The run-time increases nearly by a factor of two when
doubling the number of cores. With this setup, solving the hydrodynamics
equations needs between $40\%$ and $50\%$ of the computation time and the
radiation transport the remaining $60\%$ to $50\%$, however, these numbers are
strongly problem-dependent. Therefore even up to 1024 cores, we see good
agreement with ideal scaling. According to Amdahl’s law the full code,
including the original part of PLUTO and our implementation of the radiation
transport, is well parallelised.
## 5 Summary and conclusions
We described the implementation of a new radiation module to the PLUTO code.
The module solves for the flux-limited diffusion approximation in the two-
temperature approach. For discretisation the finite volume method is used, and
the resulting difference equations couple the updates of the temperature and
radiation energy density. Due to possibly severe time step limitations, the
set of equations is solved implicitly. For treating the non-linearity of the
temperature in the matter-radiation coupling term, we utilize the method of
Commerçon et al. (2011).
The accuracy of the implementation has been verified using different physical
and numerical setups. The first set of tests deals with purely radiative
problems that include the purely diffusive evolution towards an equilibrium,
and special setups to test the coupling terms between radiative and thermal
energy. A newly developed setup checks for the correct inclusion of the
irradiation from a central source in a spherical coordinate system.
In the second test suite we study the full simultaneous evolution of
hydrodynamics and radiation. First, sub- and super-critical radiative shock
simulations are performed and their outcomes agree very well with published
results of identical setups. Finally, we study the onset of convection in
internally heated viscous discs, and find very good agreement between 3
different, independent hydrodynamical codes. This last test also allowed us to
test the correct implementation in a spherical coordinate system and a non-
equidistant logarithmic grid. Our numerical performance tests indicate
excellent parallel scaling, up to at least 1024 processors.
The current version of the radiation module comes with routines for the
Rosseland mean opacity from Lin & Papaloizou (1985) and Bell & Lin (1994).
Additionally it is possible to use the Rosseland and Planck mean opacities
from Semenov et al. (2003).
The described radiation module can be easily used within the PLUTO
-environment. It can be found on the webpage 333http://www.tat.physik.uni-
tuebingen.de/~pluto/pluto_radiation/ as a patch for the version 4.0 of PLUTO .
###### Acknowledgements.
We gratefully thank the bwGRiD project444bwGRiD (http://www.bw-grid.de),
member of the German D-Grid initiative, funded by the Ministry for Education
and Research (Bundesministerium für Bildung und Forschung) and the Ministry
for Science, Research and Arts Baden-Württemberg (Ministerium für
Wissenschaft, Forschung und Kunst Baden-Württemberg). for the computational
resources. We gratefully acknowledges support through the German Research
Foundation (DFG) through grant KL 650/11 within the Collaborative Research
Group FOR 759: The formation of Planets: The Critical First Growth Phase. We
thank Rolf Kuiper for many stimulating discussions, either physical or
technical.
## References
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* Bitsch et al. (2013a) Bitsch, B., Boley, A., & Kley, W. 2013a, A&A, 550, A52
* Bitsch et al. (2013b) Bitsch, B., Crida, A., Morbidelli, A., Kley, W., & Dobbs-Dixon, I. 2013b, A&A, 549, A124
* Commerçon et al. (2011) Commerçon, B., Teyssier, R., Audit, E., Hennebelle, P., & Chabrier, G. 2011, A&A, 529, A35
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* Kley (1989) Kley, W. 1989, A&A, 208, 98
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* Kuiper et al. (2012) Kuiper, R., Klahr, H., Beuther, H., & Henning, T. 2012, A&A, 537, A122
* Kuiper et al. (2010) Kuiper, R., Klahr, H., Dullemond, C., Kley, W., & Henning, T. 2010, A&A, 511, A81
* Kuiper & Klessen (2013) Kuiper, R. & Klessen, R. S. 2013, A&A, 555, A7
* Levermore & Pomraning (1981) Levermore, C. D. & Pomraning, G. C. 1981, ApJ, 248, 321
* Lin & Papaloizou (1985) Lin, D. N. C. & Papaloizou, J. 1985, in Protostars and Planets II, ed. D. C. Black & M. S. Matthews, 981–1072
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|
arxiv-papers
| 2013-09-20T10:47:43 |
2024-09-04T02:49:51.207235
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Stefan M. Kolb and Matthias Stute and Wilhelm Kley and Andrea Mignone",
"submitter": "Willy Kley",
"url": "https://arxiv.org/abs/1309.5231"
}
|
1309.5260
|
# Quasi Regular Semilattice and Association Schemes in Singular Linear Space
Zhang Baohuan1 Yue Mengtian2 Li Zengti1
1.Math. and Inf. College, Langfang Teachers’ College, Langfang, 065000, China
2.Department of Scientific Research, Langfang Teachers’ College, Langfang,
065000, China [email protected]
Abstract Let $\mathbb{F}_{q}^{n+l}$ denote the $(n+l)$-dimensional singular
linear space over a finite field $\mathbb{F}_{q}$. For a fixed integer
$m\leq\min\\{n,l\\}$, denote by $\mathcal{L}^{m}_{o}(\mathbb{F}_{q}^{n+l})$
the set of all subspaces of type $(t,t_{1})$, where $t_{1}\leq t\leq m$.
Partially ordered by ordinary inclusion, one family of quasi regular
semilattices is obtained. Moreover, we obtain a association schemes and
discuss the bound of a $M$-clique.
AMS classification : 20G40, 05B35
Key words: Quasi regular semilattice; Association scheme; $M$-clique
## 1 Introduction
It is well known that lattice is an important part of poset’s theory. its
theory play an important role in many branches of mathematics, such as
computer logical design. The results on the lattices generated by transitive
sets of subspaces under finite classical groups may be found in Huo, Liu and
Wan [8, 9, 10]. In [4], Guo discussed the lattices associated with finite
vector spaces and finite affine spaces. In [12], P. Delsarte discussed the
regular semilattices in finite vector spaces. In this paper, we obtain a new
quasi regular semilattice and a new association scheme in the singular linear
space. More over we discuss the bounds of a $M$-clique.
The rest of this paper is organized as followed. In section 2, we discuss some
definitions and terminologies about lattices, regular semilattices,
association schemes and $M$-cliques. In section 3, we construct a family of
quasi regular semilattice, and then compute its parameters. In section 4, we
obtain a new association. In section 5, we discuss the bound of a $M$-clique.
## 2 Preliminaries
Let $(P,\leq)$ be a poset. We write $a<b$ whenever $a\leq b$ and $a\neq b$. If
$P$ has the minimum (respectively maximum) element, then we denote it by 0
(respectively $\upharpoonleft$), and say that $P$ is a poset with 0
(respectively $\upharpoonleft$). A poset $P$ is said to be a semilattice if
$a\wedge b:=\rm{inf}\\{a,b\\}$ exist for any two elements $a,b\in P$. Let P be
a finite poset with 0. If there is a function $r$ from $P$ to set of all the
nonnegative integers such that
* (1)
r(0)=0,
* (2)
$r(b)=r(a)+1$, if $a\lessdot b$.
Then $r$ is said to be the rank function on $P$. Note that the rank function
on $P$ is unique if it exists.
Let $P$ be a semilattice, and let $P=X_{0}\cup X_{1}\cup\cdots\cup X_{m}$,
where $X_{i}=\\{x\in P|r(x)=i\\},i=0,1,\cdots,m.$ The semilattice $(P,\leq)$
is called regular if the following three properties hold:
* (i)
Given $y\in X_{m},z\in X_{r}$ with $z\leq y$, the number of points $u\in
X_{s}$ such that $z\leq u\leq y$ is a constant $\mu(r,s)$.
* (ii)
Given $u\in X_{s}$ , the number of points $z\in X_{r}$, such that $z\leq u$ is
a constant $\nu(r,s)$.
* (iii)
Given $a\in X_{r},,y\in X_{m}$, with $a\wedge y\in X_{j}$, the number of pairs
$(b,z)\in X_{s}\times X_{m}$ such that $b\leq z,b\leq y,a\leq z$ is a constant
$\pi(j,r,s).$
In this paper, we define the concept of quasi regular semilattice as follows.
The semilattice $(P,\leq)$ is called quasi regular if the following three
properties hold:
* (i)
Given $y\in X^{m^{\prime}}_{m},z\in X^{r^{\prime}}_{r}$ with $z\leq y$, the
number of points $u\in X^{s^{\prime}}_{s}$ such that $z\leq u\leq y$ is a
constant $\mu(r(r^{\prime}),s(s^{\prime});m^{\prime})$.
* (ii)
Given $u\in X^{s^{\prime}}_{s}$ , the number of points $z\in
X^{r^{\prime}}_{r}$, such that $z\leq u$ is a constant
$\nu(r(r^{\prime}),s(s^{\prime}))$.
* (iii)
Given $a\in X^{r^{\prime}}_{r},y\in X^{m_{1}}_{m}$, with $a\wedge y\in
X^{j^{\prime}}_{j}$, the number of pairs $(b,z)\in X^{s^{\prime}}_{s}\times
X^{m^{\prime}}_{m}$ such that $b\leq z,b\leq y,a\leq z$ is a constant
$\pi(j(j^{\prime}),r(r^{\prime}),s(s^{\prime});m^{\prime}).$
Here $X_{i}=X_{i}^{0}\cup X_{i}^{1}\cup\cdots\cup X_{i}^{i}$, and
$X_{i}^{j}\cap X_{i}^{k}=\emptyset$ for $j\neq k.$
Let $X$ be a finite set of vertices. A $d$-class association scheme on $X$
consists of a set of $d+1$ symmetric relations $R_{0},R_{1},\cdots,R_{d}$ on
$V$, with identity relation $R_{0}=\\{(x,x)|x\in V\\},$ such that any two
vertices are in precisely one relation. Denoted by $(X,\\{R_{i}\\}_{0\leq
i\leq d})$. Furthermore, there are intersection numbers $p^{k}_{ij}$ such that
for any $(x,y)\in R_{k}$, the number of vertices $z$ such that $(x,z)\in
R_{i}$ and $(z,y)\in R_{j}$ equals $p^{k}_{ij}$.
The nontrivial relations can be considered as graphs, which in our case are
undirected. One immediately sees that the respective graphs are regular with
degree $v_{i}=p_{ii}^{0}$. For the corresponding adjacency matrices $A_{i}$
the axioms of the scheme are equivalent to
$\sum_{i=0}^{d}A_{i}=J,\ A_{0}=I,\ A_{i}=A_{i}^{T},\
A_{i}A_{j}=\sum_{k=0}^{d}p_{ij}^{k}A_{k}.$
It follows that the adjacency matrices generate a $(d+1)$-dimensional
commutative algebra $\mathfrak{A}$ of symmetric matrices. This algebra was
first studied by Bose and Mesner [19] and is called the Bose-Mesner algebra of
the scheme.
A nonzero vector $\alpha\in\mathbb{R}^{|X|}$ is said to be $\theta$-positive ,
if $C_{i}(\alpha)=v_{i}^{-1}<\alpha,A_{i}\alpha>\geq 0$ holds for every $i$.
Let $M$ be a subset of $[0,n]$, containing $0$, Any positive vector $\alpha$
is called an $M$-clique if it satisfies
$C_{j}(\alpha)=v_{j}^{-1}<\alpha,A_{j}\alpha>=0,\ for\ all\ j\not\in M$
Let $x_{0}$ be a point of $X$ and $s$ be an integer in $[0,n]$, A given
$\theta$-positive vector $\alpha\in\mathbb{R}^{|X|}$ is called a unicoloured
vector of center $x_{0}$ and colour $s$, if the following condition holds:
$\alpha(x)=0,\ unless\ (x,x_{0})\in R_{s}$
Let $\mathbb{F}_{q}$ be a finite field with $q$ elements, where $q$ is a prime
power. For two non-negative integers $n$ and $l$, $\mathbb{F}_{q}^{n+l}$
denotes the $(n+l)$-dimensional row vector space over $\mathbb{F}_{q}$. The
set of all $(n+l)\times(n+l)$ nonsingular matrices over $\mathbb{F}_{q}$ of
the form
$\left(\begin{array}[]{cc}T_{11}&T_{12}\\\ 0&T_{22}\end{array}\right),$
where $T_{11}$ and $T_{22}$ are nonsingular $n\times n$ and $l\times l$
matrices, respectively, forms a group under matrix multiplication, called the
singular general linear group of degree $n+l$ over $\mathbb{F}_{q}$ and
denoted by $GL_{n+l,n}(\mathbb{F}_{q})$.
Let $P$ be an $m$-dimensional subspace of $\mathbb{F}_{q}^{n+l}$, denote also
by $P$ a $m\times(n+l)$ matrix of rank $m$ whose rows span the subspace $P$
and call the matrix $P$ a matrix representation of the subspace $P$ . There is
an action of $GL_{n+l,n}(\mathbb{F}_{q})$ on $\mathbb{F}_{q}^{n+l}$ defined as
follows
$\begin{array}[]{rrr}\mathbb{F}_{q}^{n+l}\times
GL_{n+l,n}(\mathbb{F}_{q})&\longrightarrow&\mathbb{F}_{q}^{n+l},\\\
((x_{1},\cdots,x_{n},x_{n+1},\cdots,x_{n+l}),T)&\longmapsto&(x_{1},\cdots,x_{n},x_{n+1},\cdots,x_{n+l})T.\end{array}$
The above action induces an action on the set of subspaces of
$\mathbb{F}_{q}^{n+l}$; i.e., a subspace $P$ is carried by $T\in
GL_{n+l,n}(\mathbb{F}_{q})$ to the subspace $PT$. The vector space
$\mathbb{F}_{q}^{n+l}$ together with the above group action is called the
$(n+l)$-dimensional singular linear space over $\mathbb{F}_{q}$.
For $1\leq i\leq n+l$, let $e_{i}$ be the row vector in $\mathbb{F}_{q}^{n+l}$
whose $i$-th coordinate is $1$ and all other coordinates are $0$. Denote by
$E$ the $l$-dimensional subspace of $\mathbb{F}_{q}^{n+l}$ generated by
$e_{n+1},e_{n+2},\cdots,e_{n+l}$. A $m$-dimensional subspace $P$ of
$\mathbb{F}_{q}^{n+l}$ is called a subspace of type $(m,k)$ if $\dim(P\cap
E)=k$.
For a fixed subspace $P$ of type $(m_{1},k_{1})$ in $\mathbb{F}_{q}^{n+l}$,
let $\mathcal{M}^{\prime}(m_{1},k_{1};m,k;n+l,n)$ denote the set of all the
subspaces of type $(m,k)$ containing $P$.
$\mathcal{M}^{\prime}(m_{1},k_{1};m,k;n+l,n)$ is non-empty if and only if
$0\leq k_{1}\leq k\leq l$ and $0\leq m_{1}-k_{1}\leq m-k\leq n$. Let
$N^{\prime}(m_{1},k_{1};m,k;n+l,n)=|\mathcal{M}^{\prime}(m_{1},k_{1};m,k;n+l,n)|.$
The formula is given by Kaishun Wang, Jun Guo, Fenggao Li (see [7]) as
follows.
$N^{\prime}(m_{1},k_{1};m,k;n+l,n)=q^{(l-k)(m-k-
m_{1}+k_{1})}\left[n-(m_{1}-k_{1})\atop(m-k)-(m_{1}-k_{1})\right]_{q}\left[l-k_{1}\atop
k-k_{1}\right]_{q}.$
###### Lemma 2.1
Let $0\leq k_{1}\leq k_{2}\leq k\leq l,0\leq m_{1}-k_{1}\leq m_{2}-k_{2}\leq
m-k\leq n.$ For $i=1,2,3$, assume that $U_{i}$ be a set of all the subspaces
of type $(m_{i},k_{i})$ in $\mathbb{F}_{q}^{n+l}$, and $U$ be a set of all the
subspaces of type $(m,k)$ in $\mathbb{F}_{q}^{n+l}$ such that $U_{1}\subseteq
U_{2}\subseteq U$ and $U_{3}\subseteq U$. if $U_{3}\cap U_{2}=U_{1}$, then the
number of $U_{3}$ is
$q^{(\delta_{3}-\delta_{1})(\delta_{2}-\delta_{1}+k-k_{3})+(k_{3}-k_{1})(k_{2}-k_{1})}\left[\delta-\delta_{2}\atop\delta_{3}-\delta_{1}\right]_{q}\left[k-k_{2}\atop
k_{3}-k_{1}\right]_{q},$
where $\delta=m-k,\delta_{i}=m_{i}-k_{i},i=1,2,3$.
Proof. Let
$\delta_{1}=m_{1}-k_{1},\delta_{2}=m_{2}-k_{2},\delta_{3}=m_{3}-k_{3},\delta=m-k$.
By the transitivity of $GL_{n+l,n}(\mathbb{F}_{q})$ on the set of subspaces of
the same type, we may assume that
$\begin{array}[]{l}\begin{array}[]{llll}\ \ \ \ \ \ \ \ \
\delta_{1}&n-\delta_{1}&k_{1}&l-k_{1}\end{array}\\\
U_{1}=\left(\begin{array}[]{llll}I&0&0&0\\\
0&0&I&0\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
\delta_{1}\\\ k_{1}\end{array}\end{array}$
$\begin{array}[]{l}\begin{array}[]{llllll}\ \ \ \ \ \ \ \ \
\delta_{1}&\delta_{2}-\delta_{1}&n-\delta_{2}&k_{1}&k_{2}-k_{1}&l-k_{2}\end{array}\\\
U_{2}=\left(\begin{array}[]{llllll}I&0&0&0&0&0\\\ 0&0&0&I&0&0\\\
0&I&0&0&0&0\\\
0&0&0&0&I&0\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
\delta_{1}\\\ k_{1}\\\ \delta_{2}-\delta_{1}\\\
k_{2}-k_{1}\end{array}\end{array}$
and
$\begin{array}[]{l}\begin{array}[]{llllllll}\ \ \ \ \ \ \ \ \ \
\delta_{1}&\delta_{2}-\delta_{1}&\delta-\delta_{2}&n-\delta&k_{1}&k_{2}-k_{1}&k-k_{2}&l-k\end{array}\\\
U=\left(\begin{array}[]{llllllll}I&0&0&0&0&0&0&0\\\ 0&0&0&0&I&0&0&0\\\
0&I&0&0&0&0&0&0\\\ 0&0&0&0&0&I&0&0\\\ 0&0&I&0&0&0&0&0\\\ 0&0&0&0&0&0&I&0\\\
\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
\delta_{1}\\\ k_{1}\\\ \delta_{2}-\delta_{1}\\\ k_{2}-k_{1}\\\
\delta-\delta_{2}\\\ k-k_{2}\end{array}\end{array}$
Since $U_{3}\cap U_{2}=U_{1},$ we have
$\begin{array}[]{l}\begin{array}[]{llllllll}\ \ \ \ \ \ \ \ \
\delta_{1}&\delta_{2}-\delta_{1}&\delta-\delta_{2}&n-\delta&k_{1}&k_{2}-k_{1}&k-k_{2}&l-k\end{array}\\\
U_{3}=\left(\begin{array}[]{llllllll}I&0&0&0&0&0&0&0\\\ 0&0&0&0&I&0&0&0\\\
0&u_{32}&u_{33}&0&0&u_{36}&u_{37}&0\\\ 0&0&0&0&0&u_{46}&u_{47}&0\\\
\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
\delta_{1}\\\ k_{1}\\\ \delta_{3}-\delta_{1}\\\ k_{3}-k_{1}\\\
\end{array}\end{array}$
where rank $u_{47}=k_{3}-k_{1}$. Note that there are $\left[k-k_{2}\atop
k_{3}-k_{1}\right]_{q}$ choices for $u_{47}$. By the transitivity of
$GL_{n+l,n}(\mathbb{F}_{q})$ on the set of subspaces of the same type, the
number of $U_{3}$’s does not depend on the particular choice of $u_{47}$. Pick
$u_{47}=(I^{(k_{3}-k_{1})},0)$. Then $U_{3}$ has a matrix representation
$\begin{array}[]{l}\begin{array}[]{lllllllll}\ \
\delta_{1}&\delta_{2}-\delta_{1}&\delta-\delta_{2}&n-\delta&k_{1}&k_{2}-k_{1}&k_{3}-k_{1}&k-k_{2}-k_{3}+k_{1}&l-k\end{array}\\\
\left(\begin{array}[]{lllllllll}I&0&0&0&0&0&0&0&0\\\ 0&0&0&0&I&0&0&0&0\\\
0&u_{32}&u_{33}&0&0&u_{36}&0&u^{\prime}_{37}&0\\\ 0&0&0&0&0&u_{46}&I&0&0\\\
\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
\delta_{1}\\\ k_{1}\\\ \delta_{3}-\delta_{1}\\\ k_{3}-k_{1}\\\
\end{array}\end{array}$
where rank $u_{33}=\delta_{3}-\delta_{1}$. Therefore the number of $U_{3}$ is
equal to
$q^{(\delta_{3}-\delta_{1})(\delta_{2}-\delta_{1}+k-k_{3})+(k_{3}-k_{1})(k_{2}-k_{1})}\left[\delta-\delta_{2}\atop\delta_{3}-\delta_{1}\right]_{q}\left[k-k_{2}\atop
k_{3}-k_{1}\right]_{q}.$
$\hfill\Box\vskip 6.0pt plus 2.0pt minus 2.0pt$
## 3 The quasi regular semilattice
For a fixed integer $m\leq\min\\{n,l\\}$, denote by
$\mathcal{L}^{m}(\mathbb{F}_{q}^{n+l})$ the set of all subspaces of type
$(t,t_{1})$, where $t_{1}\leq t\leq m$. If we partially order
$\mathcal{L}^{m}(\mathbb{F}_{q}^{n+l})$ by the ordinary inclusion, then
$\mathcal{L}^{m}(\mathbb{F}_{q}^{n+l})$ is a semilattice, denoted by
$\mathcal{L}^{m}_{o}(\mathbb{F}_{q}^{n+l})$. For any
$A\in\mathcal{L}^{m}_{o}(\mathbb{F}_{q}^{n+l})$, The rank function of
$\mathcal{L}_{o}(\mathbb{F}_{q}^{n+l})$ is defined as follows
$r(A)=\dim(A).$
Let
$X_{i}=\\{B\in\mathcal{L}^{m}_{o}(\mathbb{F}_{q}^{n+l})|r(B)=i\\},$
and
$X_{i}^{j}=\\{B\in X_{i}|\dim(B\cap E)=j\\},j=0,1,\cdots,i,$
where $E=\langle
e_{n+1},e_{n+2},\cdots,e_{n+l}\rangle\in\mathbb{F}_{q}^{n+l}.$ We will prove
that $\mathcal{L}^{m}_{o}(\mathbb{F}_{q}^{n+l})$ is a quasi regular
semilattice and compute its parameters.
###### Lemma 3.1
Let $A\in X^{r_{1}}_{r}$, $C\in X^{m_{1}}_{m}$ and $A\leq C$. Then the number
of $B\in X^{s_{1}}_{s}$ such that $A\leq B\leq C$ is equal to
$\mu(r(r_{1}),s(s_{1});m_{1})=q^{(s-s_{1}-r+r_{1})(m_{1}-s_{1})}\left[\begin{array}[]{l}m-r+r_{1}-m_{1}\\\
s-s_{1}-r+r_{1}\end{array}\right]_{q}\left[\begin{array}[]{l}m_{1}-r_{1}\\\
s_{1}-r_{1}\end{array}\right]_{q}$
Proof. By the transitivity of $GL_{n+l}(\mathbb{F}_{q})$ on the set of
subspaces of the same type, we may assume that
$\begin{array}[]{l}\begin{array}[]{llll}\ \ \ \ \ \ \ \
r-r_{1}&n-r+r_{1}&r_{1}&l-r_{1}\end{array}\\\
A=\left(\begin{array}[]{llll}I&0&0&0\\\
0&0&I&0\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
r-r_{1}\\\ r_{1}\end{array}\end{array}$
and
$\begin{array}[]{l}\begin{array}[]{llllll}\ \ \ \ \ \ \ \
r-r_{1}&m-m_{1}-r+r_{1}&n-m+m_{1}&r_{1}&m_{1}-r_{1}&l-m_{1}\end{array}\\\
C=\left(\begin{array}[]{llllll}I&0&0&0&0&0\\\ 0&0&0&I&0&0\\\ 0&I&0&0&0&0\\\
0&0&0&0&I&0\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
r-r_{1}\\\ r_{1}\\\ m-m_{1}-r+r_{1}\\\ m_{1}-r_{1}\end{array}.\end{array}$
Since
$A\leq B\leq C,$
we have
$\begin{array}[]{l}\begin{array}[]{llllll}\ \ \ \ \ \ \ \
r-r_{1}&m-m_{1}-r+r_{1}&n-m+m_{1}&r_{1}&m_{1}-r_{1}&l-m_{1}\end{array}\\\
B=\left(\begin{array}[]{llllll}I&0&0&0&0&0\\\ 0&0&0&I&0&0\\\
0&u_{32}&0&0&u_{35}&0\\\
0&0&0&0&u_{45}&0\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
r-r_{1}\\\ r_{1}\\\ s-s_{1}-r+r_{1}\\\ s_{1}-r_{1}\end{array}\end{array},$
where rank $u_{45}=(s_{1}-r_{1})$. Note that there are $\left[m_{1}-r_{1}\atop
s_{1}-r_{1}\right]_{q}$ choices for $u_{45}$. By the transitivity of
$GL_{n+l}(\mathbb{F}_{q})$ on the set of subspaces of the same type, the
number of $B$’s does not depend on the particular choice of $u_{45}$. Pick
$u_{45}=(I^{(s_{1}-r_{1})},0)$. Then $B$ has a matrix representation
$\begin{array}[]{l}\begin{array}[]{lllllll}\ \
\delta_{2}&\delta_{1}-\delta_{2}&n-\delta_{1}&r_{1}&s_{1}-r_{1}&m_{1}-s_{1}&l-m_{1}\end{array}\\\
\left(\begin{array}[]{lllllll}I&0&0&0&0&0&0\\\ 0&0&0&I&0&0&0\\\
0&u_{32}&0&0&0&u^{\prime}_{35}&0\\\
0&0&0&0&I&0&0\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
\delta_{2}\\\ r_{1}\\\ \delta_{3}-\delta_{2}\\\
s_{1}-r_{1}\end{array}.\end{array}$
where $\delta_{1}=m-m_{1},\delta_{2}=r-r_{1},\delta_{3}=s-s_{1}$. Therefore
the number of $B$ is equal to
$q^{(s-s_{1}-r+r_{1})(m_{1}-s_{1})}\left[\begin{array}[]{l}m-r+r_{1}-m_{1}\\\
s-s_{1}-r+r_{1}\end{array}\right]_{q}\left[\begin{array}[]{l}m_{1}-r_{1}\\\
s_{1}-r_{1}\end{array}\right]_{q}.$
$\hfill\Box\vskip 6.0pt plus 2.0pt minus 2.0pt$
###### Lemma 3.2
If $B\in X^{s_{1}}_{s}$, then the number of $A\in X^{r_{1}}_{r}$ such that
$A\leq B$ is equal to
$\nu(r(r_{1}),s(s_{1}))=q^{(r-r_{1})(s_{1}-r_{1})}\left[s-s_{1}\atop
r-r_{1}\right]_{q}\left[s_{1}\atop r_{1}\right]_{q}.$
Proof. By the transitivity of $G_{n+l,l}(\mathbb{F}_{q})$ on the set of
subspaces of the same type, we may assume that
$\begin{array}[]{l}\begin{array}[]{llll}\ \ \ \ \ \ \ \
s-s_{1}&n-s+s_{1}&s_{1}&l-s_{1}\end{array}\\\
B=\left(\begin{array}[]{llll}I&0&0&0\\\
0&0&I&0\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
s-s_{1}\\\ s_{1}\end{array}\end{array}.$
Since
$A\leq B,$
we have
$\begin{array}[]{l}\begin{array}[]{llll}\ \ \ \ \ \ \ \
s-s_{1}&n-s+s_{1}&s_{1}&l-s_{1}\end{array}\\\
A=\left(\begin{array}[]{llll}u_{11}&0&u_{13}&0\\\
0&0&u_{23}&0\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
r-r_{1}\\\ r_{1}\end{array}\end{array},$
where rank $u_{23}=r_{1}$. Note that there are $\left[s_{1}\atop
r_{1}\right]_{q}$ choices for $u_{23}$. By the transitivity of
$GL_{n+l}(\mathbb{F}_{q})$ on the set of subspaces of the same type, the
number of $A$’s does not depend on the particular choice of $u_{23}$. Pick
$u_{23}=(I^{(r_{1})}\ 0)$. Then $A$ has a matrix representation
$\begin{array}[]{l}\begin{array}[]{lllll}s-s_{1}&n-s+s_{1}&r_{1}&s_{1}-r_{1}&l-s_{1}\end{array}\\\
\left(\begin{array}[]{lllll}u_{11}&0&0&u^{\prime}_{13}&0\\\
0&0&I&0&0\end{array}\right)\end{array}\begin{array}[]{l}\begin{array}[]{l}\\\
r-r_{1}\\\ r_{1}\end{array}\end{array}.$
Therefore the number of subspace $A$ is equal to
$\nu(r(r_{1}),s(s_{1}))=q^{(r-r_{1})(s_{1}-r_{1})}\left[s-s_{1}\atop
r-r_{1}\right]_{q}\left[s_{1}\atop r_{1}\right]_{q}.$
$\hfill\Box\vskip 6.0pt plus 2.0pt minus 2.0pt$
###### Lemma 3.3
Let $A\in X^{r_{1}}_{r}$ and $B\in X^{m_{2}}_{m}$. Assume that $A\wedge B\in
X^{j_{1}}_{j}$, and $(C,D)\in X_{s}^{s_{1}}\times X_{m}^{m_{1}}$ . If $C\leq
D,C\leq B,A\leq D$, then the number of $(C,D)$ is equal to
$\begin{array}[]{ll}&\pi(j(j_{1}),r(r_{1}),s(s_{1});m_{1})\\\ =&\sum_{0\leq
i\leq j,0\leq
i_{1}\leq\min\\{i,j_{1}\\}}q^{(s-s_{1}-i+i_{1})(j-j_{1}-i+i_{1}+m_{1}-s_{1})+(s_{1}-i_{1})(j_{1}-i_{1})}\\\
&\times\left[(m-m_{1})-(j-j_{1})\atop(s-s_{1})-(i-i_{1})\right]_{q}\left[m_{1}-j_{1}\atop
s_{1}-i_{1}\right]_{q}N^{\prime}(r+s-i,r_{1}+s_{1}-i_{1};m,m_{1};n+l,n).\end{array}$
Proof. Since $A\leq D,C\leq B$, we have
$A\wedge C\leq A\wedge B\leq D.$
Since
$C\leq D,$
and
$C\wedge(A\wedge B)=A\wedge(B\wedge C)=A\wedge C\in X^{i_{1}}_{i},0\leq i\leq
j,0\leq i_{1}\leq\min\\{i,j_{1}\\}.$
For $i\in[0,j]$ and $i_{1}\in[0,\min\\{i,j_{1}\\}]$, by Lemma 2.1, the number
of $C$ is equal to
$q^{(s-s_{1}-i+i_{1})(j-j_{1}-i+i_{1}+m_{1}-s_{1})+(s_{1}-i_{1})(j_{1}-i_{1})}\left[(m-m_{1})-(j-j_{1})\atop(s-s_{1})-(i-i_{1})\right]_{q}\left[m_{1}-j_{1}\atop
s_{1}-i_{1}\right]_{q}.$
Since $C\leq D,A\leq D$, $C+A\leq D$. It follows from $A\wedge C\in
X_{i}^{i_{1}}$ that $A+C$ is a subspace of type $(r+s-i,r_{1}+s_{1}-i_{1})$.
Therefore, the number of $D$ is equal to
$N^{\prime}(r+s-i,r_{1}+s_{1}-i_{1};m,m_{1};n+l,n)$. Hence the desired result
follows. $\hfill\Box\vskip 6.0pt plus 2.0pt minus 2.0pt$
By Lemma 3.1,Lemma 3.2 and Lemma 3.3, we obtain the following theorem.
###### Theorem 3.4
Semilattice $\mathcal{L}_{0}^{m}(\mathbb{F}_{q}^{n+l})$ is a quasi regular
semilattice. Its parameters are given by the formulas
$\mu(r(r_{1}),s(s_{1});m_{1})=q^{(s-s_{1}-r+r_{1})(m_{1}-s_{1})}\left[\begin{array}[]{l}m-r+r_{1}-m_{1}\\\
s-s_{1}-r+r_{1}\end{array}\right]_{q}\left[\begin{array}[]{l}m_{1}-r_{1}\\\
s_{1}-r_{1}\end{array}\right]_{q},$
$\nu(r(r_{1}),s(s_{1}))=q^{(r-r_{1})(s_{1}-r_{1})}\left[s-s_{1}\atop
r-r_{1}\right]_{q}\left[s_{1}\atop r_{1}\right]_{q},$
and
$\begin{array}[]{ll}&\pi(j(j_{1}),r(r_{1}),s(s_{1});m_{1})\\\ =&\sum_{0\leq
i\leq j,0\leq i_{1}\leq
i}q^{(s-s_{1}-i+i_{1})(j-j_{1}-i+i_{1}+m_{1}-s_{1})+(s_{1}-i_{1})(j_{1}-i_{1})}\\\
&\times\left[(m-m_{1})-(j-j_{1})\atop(s-s_{1})-(i-i_{1})\right]_{q}\left[m_{1}-j_{1}\atop
s_{1}-i_{1}\right]_{q}N^{\prime}(r+s-i,r_{1}+s_{1}-i_{1};m,m_{1};n+l,n).\end{array}$
## 4 Association schemes
Given $d$ and $d_{0}$ with $0\leq d_{0}\leq d\leq\min\\{n,l\\}$. For
$h=\min\\{k,d_{0}\\},$ we define
$\widetilde{X_{k}}=\left\\{\begin{array}[]{ll}X_{k}^{0}\cup\cdots\cup
X_{k}^{h},&if\ 0\leq k\leq d-1,\\\ X_{d}^{d_{0}},&if\ k=d,\end{array}\right.$
and
$R_{k}=\\{(B,C)\in X_{d}^{d_{0}}\times X_{d}^{d_{0}}|B\wedge
C\in\widetilde{X_{d-k}}\\}.$
Then the $R=\\{R_{0},R_{1},\cdots,R_{d}\\}$ is a partition of
$X_{d}^{d_{0}}\times X_{d}^{d_{0}}$. For any $B,C\in X_{d}^{d_{0}}$, let
$A_{k(k^{\prime})}(B,C)=\left\\{\begin{array}[]{ll}1,&if\ B\wedge C\in
X^{k^{\prime}}_{d-k},\\\ 0,&otherwise.\end{array}\right.$
Then, for $h=\min\\{d-k,d_{0}\\}$, we define
$A_{k}=\left\\{\begin{array}[]{ll}A_{k(0)}+A_{k(1)}+\cdots+A_{k(h)},&if\ 1\leq
k\leq d,\\\ A_{0(d_{0})},&if\ k=0.\end{array}\right.$
Let
$D_{i,k}(B,C)=\left\\{\begin{array}[]{ll}1,&if\ B\leq C,\\\
0,&otherwise.\end{array}\right.$
The $D_{i,k}$ is a Riemann matrix indexed by
$\widetilde{X_{i}}\times\widetilde{X_{k}}$. For convenience, we write
$D_{i}=D_{i,d}$. Let $\mathfrak{A}=\langle A_{0},A_{2},\cdots,A_{d}\rangle$
denote the $(d+1)$-dimensional real vector space generated by the
$A_{0},A_{2},\cdots,A_{d}$. Let $C_{i}=D^{T}_{i}D_{i}$, for $i=0,1,\cdots,d.$
We have the lemma as follows.
###### Lemma 4.1
The matrices $C_{0},C_{1},\cdots,C_{d}$ generate $\mathfrak{A}$, with
$C_{t}=\sum_{h=0}^{\min\\{t,d_{0}\\}}\nu(t(h),d(d_{0}))A_{0}+\sum_{k=1}^{d}\sum_{\lambda=0}^{\min\\{d-k,d_{0}\\}}\sum_{h=0}^{\min\\{t,d_{0}\\}}\nu(t(h),(d-k)(\lambda))A_{k},$
(1)
where $t=0,1,\cdots,d.$
Proof. For any $(A,B)\in X_{d}^{d_{0}}\times X_{d}^{d_{0}}$, the number of
$F\in\widetilde{X_{t}}$ such that $F\leq A\wedge B$ is the $(A,B)$-entry in
two members of (1). It implies that the system (1) holds. On the other hand,
its matrix has rank $d+1$. Hence, the result is obtained. $\hfill\Box\vskip
6.0pt plus 2.0pt minus 2.0pt$
###### Lemma 4.2
Let $A\in X^{d_{1}}_{d},B\in X^{s_{1}}_{s}$, then the number of $D\in
X^{r_{1}}_{r}$ such that $D\leq A$ and $B\wedge D\in X^{j_{1}}_{j}$ is a
constant $\psi(j(j_{1}),r(r_{1}),s(s_{1});d_{1})$.
Proof. For a given $k$ and $k_{1}$, with $0\leq k\leq h=\min\\{r,s\\}$ and
$0\leq k_{1}\leq k$. We assume that $(C,D)\in(X_{k}^{k_{1}},X_{r}^{r_{1}})$
and $C\leq B,C\leq D\leq A.$ Counting the number of the pairs $(C,D)$ in two
different ways, we obtain
$\sum_{j=0}^{h}\sum_{j_{1}=0}^{j}\nu(k(k_{1}),j(j_{1}))\psi(j(j_{1}),r(r_{1}),s(s_{1});d_{1})=\nu(k(k_{1}),s(s_{1}))\mu(k(k_{1}),r(r_{1});d_{1}).$
(2)
This yields a system $(\ref{el2})$ of linear equations in unknown
$\psi(j(j_{1}),r(r_{1}),s(s_{1});d_{1})$ with fixed $r(r_{1}),s(s_{1})$. Since
the matrix $[\nu(k(k_{1}),j(j_{1}))]$ of the system $(\ref{el2})$ is a
nonsingular, these equations uniquely determine the
$\psi(j(j_{1}),r(r_{1}),s(s_{1});d_{1})$. $\hfill\Box\vskip 6.0pt plus 2.0pt
minus 2.0pt$
###### Theorem 4.3
The system $(X_{d}^{d_{0}},R)$ is an association scheme.
Proof. We establish the matrix relation as follows
$C_{r}C_{s}=\sum_{k=0}^{d}\sum_{\lambda=0}^{h_{1}}\left\\{\sum_{j=0}^{r}\sum_{j^{\prime}=0}^{h_{2}}\left[\sum_{r^{\prime}=0}^{h_{3}}\left(\psi(j(j^{\prime}),r(r^{\prime}),(d-k)(\lambda);d_{0})\sum_{s^{\prime}=0}^{h_{4}}\pi(j(j^{\prime}),r(r^{\prime}),s(s^{\prime});d_{0})\right)\right]\right\\}A_{k}.$
(3)
Here
$h_{1}=\min\\{d-k,d_{0}\\},h_{2}=\min\\{j,d_{0}\\},h_{3}=\min\\{r,d_{0}\\},h_{4}=\min\\{s,d_{0}\\}.$
In fact, By the definition, for any $D,F\in X_{d}^{d_{0}}$, the $(D,F)$-entry
of $C_{r}C_{s}$ is the number of triples
$(A,B,C)\in\widetilde{X_{r}}\times\widetilde{X_{s}}\times\widetilde{X_{d}}$
such that $A\leq D\wedge C$ and $B\leq F\wedge C$ hold. Let $D\wedge F\in
X_{d-k}^{\lambda}$. We assume that $A\wedge F\in X_{j}^{j^{\prime}}$ with
fixed $j$ and $j^{\prime}$, by Lemma 4.2, there are
$\sum_{r^{\prime}=0}^{\min\\{r,d_{0}\\}}\psi(j(j^{\prime}),r(r^{\prime}),(d-k)(\lambda);d_{0})$
choices for $A\leq D$. For each $A$, by Lemma 3.3, there are
$\sum_{s^{\prime}=0}^{\min\\{s,d_{0}\\}}\pi(j(j^{\prime}),r(r^{\prime}),s(s^{\prime});d_{0})$
choices for pairs $(B,C)$ with $B\leq F\wedge C,A\leq C$. Hence, (3) holds. By
Lemma 4.1, $\mathfrak{A}$ is a Bose-Mesner algebra. It implies the assertion.
$\hfill\Box\vskip 6.0pt plus 2.0pt minus 2.0pt$
## 5 $M$-clique
In this section, we always assume that $l=0$, $X=X_{d}^{0}$. By Theorem 4.3,
we obtain a association scheme $\mathfrak{X}=(X,\\{R_{i}\\}_{0\leq i\leq d})$
in $\mathbb{F}_{q}^{n}$. Next, we will discuss the bound of $M$-clique on
association scheme $(X,\\{R_{i}\\}_{0\leq i\leq d})$.
We associate the following linear programming problem, in the real variables
$\xi_{k}(k\in K=\\{k\in[0,d]|m_{k}>1\\}):$(see [5])
$\left\\{\begin{array}[]{l}{\rm min}(f^{\prime})=1+\sum_{k\in
K}(1-b_{k})\xi_{k},\\\ \hbox{subject to}\ \xi_{k}\geq 0,\ \hbox{all}\ k\in
K,\\\ \hbox{and to}\ \sum_{k\in K}(P_{i}(k)-b_{k}v_{i})\xi_{k}\leq v_{i},\
\hbox{all}\ i\in M\backslash\\{0\\}.\end{array}\right.$ (4)
Here $P_{i}(k)$ is eigenvalue of a association scheme, $v_{i}$ is degree of a
association scheme, $b_{k}=(v_{s}^{-1}P_{s}(k))^{2},k\in[0,d]$.
###### Proposition 5.1
([5]) The problem (4) admits a solution $\min(f^{\prime})$, and any
unicoloured $M$-clique $\alpha$ satisfies
$\hat{c}_{0}(\alpha)\leq\min(f^{\prime})$.
By Proposition 5.1, we obtain the following theorem.
###### Theorem 5.2
Let $\alpha\in\mathbb{R}^{|X|}$ be a unicoloured $M$-clique on association
scheme $\mathfrak{X}=(X,\\{R_{i}\\}_{0\leq i\leq d})$ with $M=\\{0,1\\}$ and
colour $s=1$. then
$\hat{c}_{0}(\alpha)=<\alpha,\alpha>^{-1}<\alpha,1>^{2}\leq\frac{q^{d+1+n}-3q^{2d+1}+5q^{d+1}-2q^{d}+2q^{n+d}-3q^{n+1}}{q^{n+d}-q^{2d+1}+2q^{d+1}-q^{d}-q^{n+1}}.$
Proof. By the ([12]), the eigenvalues of an association scheme of
$\mathfrak{X}=(X,\\{R_{i}\\}_{0\leq i\leq d})$ are given by the formula
$P_{i}(r)=\sum_{j=0}^{i}(-1)^{i-j}\left[d-j\atop i-j\right]_{q}\left[d-r\atop
j\right]_{q}\left[n-d+j-r\atop j\right]_{q}q^{rj+{i-j\choose 2}}.$
Pick $\xi_{k}=\xi\delta_{k,1}$. It is known that the degree of the
$\mathfrak{X}=(X,\\{R_{i}\\}_{0\leq i\leq d})$ are given by the formula
$v_{i}=q^{i^{2}}\left[n-d\atop i\right]_{q}\left[d\atop i\right]_{q}.$
By the above formulas, we obtain
$v_{1}=\frac{q(q^{n-d}-1)(q^{d}-1)}{(q-1)^{2}},$
$P_{1}(1)=-\frac{q^{d}-1}{q-1}+\frac{q(q^{d-1}-1)(q^{n-d}-1)}{(q-1)^{2}},$
For $s=1,i=1$, we have $(P_{1}(1)-b_{1}v_{1})\xi\leq v_{1}$,
$(P_{1}(1)-(v_{1}^{-1}P_{1}(1))^{2}v_{1})\xi\leq v_{1}$. It implies that
$(v_{1}P_{1}(1)-P_{1}(1)^{2})\xi\leq v_{1}^{2}$. Since
$v_{1}-P_{1}(1)=\frac{q^{n}-1}{q-1}>0.$
We have
$\xi\leq\frac{v_{1}^{2}}{p_{1}(1)(v_{1}-p_{1}(1)}.$
and
$1-b_{1}=1-\frac{p_{1}(1)^{2}}{v_{1}^{2}}=\frac{v_{1}^{2}-p_{1}(1)^{2}}{v_{1}^{2}}>0.$
It follows that
$\begin{array}[]{ll}\min(f^{\prime})&=1+(1-b_{1})\xi\\\ &\leq
1+\frac{v_{1}+p_{1}(1)}{p_{1}(1)}\\\
&=\frac{q^{d+1+n}-3q^{2d+1}+5q^{d+1}-2q^{d}+2q^{n+d}-3q^{n+1}}{q^{n+d}-q^{2d+1}+2q^{d+1}-q^{d}-q^{n+1}}.\end{array}$
By Proposition 5.1, we have
$\hat{c}_{0}(\alpha)\leq\frac{q^{d+1+n}-3q^{2d+1}+5q^{d+1}-2q^{d}+2q^{n+d}-3q^{n+1}}{q^{n+d}-q^{2d+1}+2q^{d+1}-q^{d}-q^{n+1}}.$
$\hfill\Box\vskip 6.0pt plus 2.0pt minus 2.0pt$
## Acknowledgement
This research is supported by NSF of Hebei Province (A2013408009), NSF of
Hebei Education Department(ZH2012082), the Specialized Research Fund for the
Doctoral Program of Higher Education of China (No.20121303110005) and the
foundation of Langfang Teachers’ College (LSBS201205).
## References
* [1]
* [2] M. Aigner, Combinatorial Theory, Springer-Verlag, Berlin, 1979.
* [3] Z. Wan, Y. Huo, Lattices generated by orbits of subapaces under finite classical groups, second ed., Science Press, Beijing, 2002(in Chinese).
* [4] J. Guo, Lattices associated with finite vector spaces and finite affine spaces, Ars Combin. 88 (2008) 47–53.
* [5] P. Delsarte, Pairs of vectors in the space of an association scheme, Philips Res. Rep. 32 (1977), 373-411.
* [6] J. Guo, S.G. Gao, K. Wang, Lattices generated by subgraphs in d-bounded distance-regular graphs, Discrete mathematics, 308 (2008),5260-5264.
* [7] Kaishun Wang, Jun Guo, Fenggao Li, Singular linear space and its applications. Finite Fields and Their Applications 17(2011) 395-406.
* [8] Y. Huo, Y. Liu and Z. Wan, Lattices generated by transitive sets of subspaces under finite classical groups I, Comm. Algebra 20 (1992) 1123–1144.
* [9] Y. Huo, Y. Liu and Z. Wan, Lattices generated by transitive sets of subspaces under finite classical groups II, the orthogonal case of odd characteristic, Comm. Algebra 20 (1993) 2685–2727.
* [10] Y. Huo, Y. Liu and Z. Wan, Lattices generated by transitive sets of subspaces under finite classical groups, the orthogonal case of even characteristic III, Comm. Algebra 21 (1993a) 2351–2393.
* [11] Y. Huo and Z. Wan, On the geomericity of lattices generated by orbits of subspaces under finite classical groups, J. Algebra 243 (2001) 339–359.
* [12] P. Delsarte, Association Schemes and $t$-designs in Regular Semilattices. Journal of Combinatorial Theory(A) 20, 230-243(1976).
* [13] J. Nan and J. Guo, Lattices generated by two orbits of subspaces under finite singular classical groups, Comm. Algebra, 38 (2010) 2026–2036.
* [14] K. Wang and Y. Feng, Lattices generated by orbits of flats under finite affine groups, Comm. Algebra 34 (2006) 1691–1697.
* [15] K. Wang and J. Guo, Lattices generated by orbits of totally isotropic flats under finite affine-classical groups, Finite Fields Appl. 14 (2008) 571–578.
* [16] K. Wang and J. Guo, Lattices generated by two orbits of subspaces under finite classical groups, Finite Fields Appl. 15 (2009) 236–245
* [17] K. Wang and Z. Li, Lattices associated with vector space over a finite field, Linear. Alg Algebra Appl. 429 (2008) 439–446.
* [18] Kaishun Wang, Jun Guo, Fenggao Li,Association schemes based on attenuated spaces, European Journal of Combinatorics 31 (2010) 297-305.
* [19] R.C. Bose and D.M.Mesner, On linear associative algebras corresponding to association schemes of partially balanced designs, Ann. Math. Statist. 30 (1959), 21 C38.
* [20] S.G. Gao, J. Guo, A construction of distance-regular graphs from subspaces in d-bounded distance-regular graphs, ARS Combinatoria 98 (2011), 135-148.
|
arxiv-papers
| 2013-09-20T13:09:48 |
2024-09-04T02:49:51.219744
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Zhang Baohuan, Yue Mengtian, Li Zengti",
"submitter": "Li Zengti",
"url": "https://arxiv.org/abs/1309.5260"
}
|
1309.5275
|
# An open dataset for research on audio
field recording archives: freefield1010
Dan Stowell and Mark D. Plumbley
Centre for Digital Music, Queen Mary University of London
[email protected]
###### Abstract
We introduce a free and open dataset of 7690 audio clips sampled from the
field-recording tag in the Freesound audio archive. The dataset is designed
for use in research related to data mining in audio archives of field
recordings / soundscapes. Audio is standardised, and audio and metadata are
Creative Commons licensed. We describe the data preparation process,
characterise the dataset descriptively, and illustrate its use through an
auto-tagging experiment.
## 1 Introduction
Digital sound archives hold vast resources of material, including speech,
music, and naturalistic and ethnographic field recordings (Ranft, 2004).
However, there are still many challenges in organising and searching in these
archives. In recent decades, research fields such as automatic speech
recognition (ASR) and music information retrieval (MIR) have developed
automatic methods for labelling and transcribing specific types of sound. Yet
even if we put speech and music to one side, we still have a large and
valuable range of recorded sound, and there has been relatively little work in
organising and searching this non-speech-non-music audio. Such research may
come under the umbrella of “computational auditory scene analysis” (CASA)
(Wang and Brown, 2006). Tasks have begun to be addressed in recent years such
as automatically labelling the type of audio scene, or automatically detecting
and labelling the events within the audio scene (Giannoulis et al., 2013, and
citations therein).
Important for research development is the existence of standard datasets that
can be independently reused by researchers. This has motivated many community
efforts in ASR and in MIR. In our own recent work we created a set of “audio
scene” recordings as part of the IEEE AASP “D-CASE” challenge, consisting of
30-second binaural recordings made by three recordists in locations around
London (Giannoulis et al., 2013). Such focussed datasets are valuable for
developing algorithms for the specific tasks considered, but they are quite
different from most audio archives, being tightly calibrated in their
production and of moderate size.
In order for datasets to be relevant to applications in sound archives, they
need to be large enough that they (a) reflect the diversity of content in
audio archives and (b) give some indication of scalability issues for analysis
and for visualisation/navigation. However, large datasets are expensive and
time-consuming to record from scratch, while most existing archives cannot be
freely redistributed due to the associated copyright and licensing terms. It
is illustrative that the creators of a recent large pop-music dataset (the
Million Song Dataset) worked around this issue by distributing not the audio,
but some pre-computed features derived from the audio (Bertin-Mahieux et al.,
2011). This enabled them to distribute a large open dataset of relevance to
pop music archives, but as a consequence it restricted the types of analysis
possible: researchers are constrained to using the specific pre-computed
features they provided.
For our purposes, a notable initiative is the Freesound
archive,111http://freesound.org/ which hosts extensive holdings of
crowdsourced audio recordings, reusable under Creative Commons and public
domain licences. Established in 2005, it holds more than 160,000 sounds from
thousands of users around the world. It contains a wide range of sound types,
including field recordings, recordings from contact mics and hydrophones, and
synthetic sounds.
Freesound is already usable for various research purposes, having open
licensing conditions and an easy-to-use application programming interface
(API). However, it is a large crowdsourced and continually-updated collection,
not a fixed dataset. Files are continually added and removed, and metadata
changed; it is too large to easily redistribute among researchers; various
licences are used, not all mutually compatible; and files are in various
formats (e.g. WAV/AIFF/ MP3/Ogg, number of channels, sample-rate, duration)
which can be inconvenient for those developing algorithms. Also, a
crowdsourced archive such as Freesound is typically more heterogeneous than
traditional archives such as the British Library Sound Archive, whose curation
involves manual attention to file formats and metadata (Ranft, 2004).
We therefore chose to compile a free and open dataset of a wide range of
sounds, standardised and curated from a fraction of the extensive holdings of
the Freesound archive. This dataset is intended to be of use to researchers
developing methods for working with field recordings in audio archives. In the
following we describe how we designed and prepared the data. We then describe
how anyone can access the data and work with it, and illustrate with an
automatic recognition experiment to infer the presence of tags.
## 2 Dataset
In order to maximise the potential usefulness of the dataset, before
preparation we considered the following design criteria:
Content: Firstly we aimed to reflect the content that a general audio archive
might collect. Given the range of archive policies, from crowdsourced to
strictly curated, we opted for a middle way, by using Freesound contributions
but only those under the field-recording tag. Initial inspection of the
various tags in Freesound determined that this tag was mostly free of
interpretation issues, unlike for example ambience which Freesound users often
use for ambient field recordings but also for synthesised atmospheric
soundtrack sounds. On the other hand, we decided not to manually curate the
collection to a specific definition of “field recording”, in part because such
curation is difficult to apply to material from unknown third parties, and
also because audio archive collections are rarely so narrowly construed.
Licensing: Freesound contains material under various licences. Most common are
the Creative Commons CC-BY licence and the “CC0” public domain dedication,
though there is a small proportion of older material under other licences such
as Creative Commons “sampling” licences which are not compatible with CC-BY.
We wished to be able to apply a single open licence to the overall dataset, so
we restricted ourselves to CC-BY and CC0 material, which means the overall
dataset can be published under CC-BY. We also needed to respect the
attribution requirements in the CC-BY source material, which we implemented by
ensuring we stored the author metadata with each file as well as a URL to link
back to the original source.
Size and duration: We aimed to produce a dataset of manageable storage size,
for ease of redistribution, yet suitably diverse. We also wished to produce
sound excerpts of a standard and relatively short duration, so that they could
be used in listening tests without risking listener fatigue, and so that
automatic tests could run efficiently. These motivations led us to settle on a
fixed ten-second duration for each excerpt.
File formats: We also aimed to use a standardised file format. Freesound
allows users to upload sounds of any sample rate, any number of channels, and
in various file formats (uncompressed and lossy-compressed). We chose not to
use audio coming from lossy compressed file formats, in case of artifacts
introduced by the codecs. We also chose to convert all downloaded sounds to
standard CD-quality mono WAV files. We considered standardising on 24-bit
and/or 96 kHz as recommended in the IASA “TC04” archiving standard (IASA
Technical Committee, 2009). However, in our experience 16-bit PCM has been
more widely compatible than 24-bit PCM: the latter is not well handled by some
older versions of Matlab and some Python audio libraries, although this
situation is improving. Also the majority of the original downloaded audio was
in 44.1 kHz. We therefore settled on 16-bit 44.1 kHz. In the crowdsourced
Freesound archive, amplitude levels are uncontrolled, which may be problematic
for listening tests. We therefore chose to amplitude-normalise each excerpt in
our dataset.
Dataset partitioning: In data mining and machine-learning experiments, it is
useful to have a dataset partitioned into separate subsets—e.g. one for
training and another for testing (Witten and Frank, 2005, Chapter 5). To
facilitate this we chose to partition our data into ten equally-sized subsets.
Partitioning can be purely random or can be stratified: for example, if the
data was intended for an experiment detecting the presence of geotags, then
the partitions could be arranged such that each partition had an equal mix of
geotagged and non-geotagged data (Witten and Frank, 2005, Chapter 5). However,
our dataset is intended for various purposes and not for a single specific
experiment, so we opted for the simple random partitioning. Researchers can
choose to use these subsets for comparability with others, or to perform their
own partitioning.
### 2.1 Preparation
We first obtained specific permission from Freesound to perform our relatively
large-scale data download. We created a Python script based on the official
freesound-python code, which we used to download the files from
Freesound.222https://github.com/danstowell/freesound-python/tree/tagsearch The
script was run to download all files matching all of the following criteria:
* •
Tagged field-recording (which contained 17807 sounds in total)
* •
Length 10 seconds or greater
* •
Audio file format WAV
* •
Published under either the CC-BY licence or CC0
* •
Audio with 1 or 2 channels
* •
Audio sample format one of: pcm16, pcm24, pcm32
Each file was saved along with its metadata in JSON
format.333http://tools.ietf.org/html/rfc4627 A small number of files (44)
failed to download completely; these were detected by using the sox command-
line audio tool (v 14.3.2) to attempt WAV decoding, and deleting the files
which reported end-of-file errors or similar. The script was run in July 2013,
taking about a week to download 328 GB of material.
We then prepared a 10-second standardised excerpt from each audio file, taken
from the middle of the audio recording, and used sox to convert it to a
standardised file format: WAV, single-channel, sample rate 44.1 kHz, 16-bit
PCM, amplitude normalised to $-2$ dB (empirically selected as the maximum gain
before clipping).
We inspected the excerpts for any further issues, by listening to all
10-second extracts. We found a few (seven) which were pure silence with DC
offset. While this may sometimes be valid audio when considered in context, it
led to normalisation issues, and could be problematic for some applications;
we decided to remove these files.
The above procedure resulted in 7690 sound excerpts with accompanying
metadata. We placed the exerpts into 10 separate partitions, where the
allocation was by pseudorandom shuffle initialised with a fixed seed value for
repeatability (given in the Python script referenced above). Each of the 10
partitions has about 128 minutes of audio; the dataset totals over 21 hours of
audio.
Figure 1: Density plot of all geolocation tags in the dataset. Source code for
this plot is in Appendix A.
Each file is associated with various metadata (author, date, licence etc.)
including a median of 7 tags per file (range 1–68). Around 40% of the files
come with geolocation metadata. Figure 1 shows a density plot for all
geolocations in the dataset. It indicates a broad geographic spread, although
with much the strongest density in Europe (perhaps understandable given that
Freesound is a European project), and a relative lack of sounds tagged from
Africa and Russia.
### 2.2 Availability
The full dataset is available online, hosted on the Internet
Archive,444http://archive.org/details/freefield1010 and our institutional
repository555http://c4dm.eecs.qmul.ac.uk/rdr/handle/123456789/35, under an
overall CC-BY licence. For portability, each of the ten subsets of the data is
presented as a single zip file which can fit on a data CD (around 570 MB per
subset). We list MD5 checksums of these zip files so that their integrity can
be verified.
The metadata stored alongside the audio is easy to work with. Appendix A gives
a simple Python code example for working with the JSON metadata, in this case
the code used to generate Figure 1.
## 3 A classification experiment
To demonstrate the dataset in use, we conducted an experiment using a binary
classification paradigm using audio content analysis to infer the
presence/absence of particular tags (sometimes called auto-tagging (Ellis et
al., 2011)). For this we used the simple baseline classifier smacpy presented
in Giannoulis et al. (2013), which provides a standard implementation of a
common approach based on Gaussian mixture modelling of Mel frequency cepstral
coefficients (MFCCs).
Freesound uses a free-tagging system, with users able to associate an
arbitrary number of tags with an audio file, and also able to create new tags.
We thus manually chose a small selection of tags which appear relatively
often—birdsong, city, people, nature, train, voice, water—to use for this
experiment, creating a binary classification task for each one. We note that
we might expect the tags to have a varying directness of connection with the
audio content: we would expect recordings tagged voice to feature human voice
sounds relatively prominently, whereas recordings tagged city might contain
many sounds in ensemble, with perhaps no sonic component always present.
We further added two pseudo-tags to the study, based on other metadata
attributes: an indicator of whether or not the item comes with geolocation
data (__geotagged), and an indicator of whether the item is CC-BY licensed
(__ccby), as opposed to the only other licence present, public domain. These
pseudo-tags in principle have no direct connection to the audio content,
although there may be correlations due to circumstantial effects (for example,
geolocation might more often be stored when recording outdoor scenes). We
therefore expected only mild if any ability to predict these metadata
attributes from audio.
Table 1: Prevalence of selected tags in freefield1010. Tag | Num tagged | Proportion (%)
---|---|---
birdsong | 198 | 2.6
city | 562 | 7.3
nature | 905 | 11.8
people | 321 | 4.1
train | 411 | 5.3
voice | 556 | 7.2
water | 707 | 9.2
__geotagged | 3058 | 39.8
__ccby | 6111 | 79.5
Table 1 provides a summary of the prevalence of the selected tags in the
freefield1010dataset. It is important to note that, for the “true” tags
especially, each tag is present in only a minority of the items; the ratio of
positive to negative instances is highly skewed. This has consequences for how
we evaluate automatic classification: rather than using raw accuracy, which
fails to account for this skew, we use the area under the curve (AUC)
statistic derived from a receiver operating characteristic (ROC) curve
(Provost et al., 1998; Fawcett, 2006).
Our experiment proceeded as follows: for each of the selected tags and pseudo-
tags, we performed a ten-fold cross-validation experiment using the folds
defined by the ten subsets of freefield1010. This means that for each fold, we
used nine of the ten subsets as training data for the classifier, where the
presence/absence of the tag was the binary attribute to be learnt, and then
tested the classifier using the audio from the one remaining subset. For each
such run, we calculated the numbers of correct and incorrect decisions, and
used this to calculate the AUC statistic (as in Fawcett (2006)). Source code
for this experiment is available
online.666https://github.com/danstowell/smacpy/tree/freefield1010
Figure 2: Results for automatic inference of tag presence/ absence. Plot shows
the mean and the 95% confidence interval of the AUC score for each tag (and
pseudo-tag) studied, across ten-fold crossvalidation.
Results in Figure 2 show that the tags can be automatically inferred from
audio with varying degrees of reliability. The best result is for birdsong at
82% AUC, while the weakest result for a true tag is city at 63%. Note that the
standard interpretation of an AUC value is that it tells us the probability
that the algorithm will rank a random positive instance higher than a random
negative instance (Fawcett, 2006); chance performance is always 50% for the
AUC statistic. Results for the two pseudo-tags attain around 58% AUC: above
chance, but still very weak. It indicates that there is some mild difference
in audio content between the positive and negative instances, but much less
than for the true tags.
The classification performance for true tags appears to show some connection
with the “directness” issue raised above: the tags which yield weakest
recognition performance (city, people, nature) can be said to have an indirect
connection with the audio content. However these results are illustrative
only, using a baseline classifier rather than a leading-edge algorithm. The
strongest performance (82%) is still much lower than is desirable for a binary
classifier deployed in a live system.
The 95% confidence intervals (error bars) in Figure 2 are relatively small and
well-separated. This illustrates that the dataset is of sufficient size to
make inferences about the relative predictability of these tags from the audio
content, and also that there is relative consistency among the ten folds of
the data.
## 4 Conclusions
In this paper we have described the preparation of a free and open audio
dataset, designed primarily for use in research on data mining of audio
archives. The dataset is derived from a subset of the Freesound archive, but
fixed and standardised so as to facilitate reproducible research. In an auto-
tagging experiment we demonstrated that the dataset can be used to probe
issues such as the differential predictability of tags from audio. We hope
that the dataset will prove useful to others.
## Acknowledgments
We wish to thank the Freesound developers and maintainers at the Music
Technology Group in Universitat Pompeu Fabra, for running the excellent
Freesound archive and for allowing us API access to the data.
DS & MP are supported by an EPSRC Leadership Fellowship EP/G007144/1.
## References
* Bertin-Mahieux et al. [2011] T. Bertin-Mahieux, D. P. W. Ellis, B. Whitman, and P. Lamere. The Million Song Dataset. In _Proceedings of the 12th International Conference on Music Information Retrieval (ISMIR-11)_ , pages 591–596, Miami, FL, USA, Oct 2011.
* Ellis et al. [2011] K. Ellis, E. Coviello, and G. Lanckriet. Semantic annotation and retrieval of music using a bag of systems representation. In _Proceedings of the 12th International Conference on Music Information Retrieval (ISMIR-11)_ , pages 723–728, Miami, FL, USA, Oct 2011.
* Fawcett [2006] T. Fawcett. An introduction to ROC analysis. _Pattern recognition letters_ , 27(8):861–874, 2006. doi: 10.1016/j.patrec.2005.10.010.
* Giannoulis et al. [2013] D. Giannoulis, D. Stowell, E. Benetos, M. Rossignol, M. Lagrange, and M. D. Plumbley. A database and challenge for acoustic scene classification and event detection. In _Proceedings of the European Signal Processing Conference (EUSIPCO 2013)_ , 2013.
* IASA Technical Committee [2009] IASA Technical Committee. _Guidelines on the Production and Preservation of Digital Audio Objects_. Number IASA-TC04. International Association of Sound and Audiovisual Archives, Aarhus, Denmark, second edition, 2009. URL http://www.iasa-web.org/tc04/audio-preservation.
* Provost et al. [1998] F. J. Provost, T. Fawcett, and R. Kohavi. The case against accuracy estimation for comparing induction algorithms. In _Proceedings of the International Conference on Machine Learning (ICML)_ , pages 445–453, 1998.
* Ranft [2004] R. Ranft. Natural sound archives: Past, present and future. _Anais da Academia Brasileira de Ciências_ , 76(2):456–460, 2004. doi: 10.1590/S0001-37652004000200041.
* Wang and Brown [2006] D. L. Wang and G. J. Brown, editors. _Computational Auditory Scene Analysis: Principles, Algorithms, and Applications_. IEEE Press, New York, 2006.
* Witten and Frank [2005] I. H. Witten and E. Frank. _Data Mining: Practical Machine Learning Tools and Techniques_. Morgan Kaufmann, San Francisco, CA, USA, 2nd edition, 2005.
## Appendix A: Code example
The following Python code was used to generate Figure 1 (using Python 2.7).
⬇
import glob, json
import matplotlib.pyplot as plt
import matplotlib.cm as cm
ffpath = ’./freefield1010’
outfile = ’plots/plotgeo.pdf’
lats = [], lons = []
for onepath in glob.iglob(’%s/*/*.json’ % ffpath):
jsonfile = open(onepath, ’r’)
jsondata = json.load(jsonfile)
jsonfile.close()
geo = jsondata.get(u’geotag’)
if geo != None:
lats.append(float(geo[u’lat’]))
lons.append(float(geo[u’lon’]))
# plot
plt.figure()
plt.hexbin(lons, lats, gridsize=25, bins=’log’, cmap=cm.binary)
plt.xlabel(’Longitude’)
plt.ylabel(’Latitude’)
cb = plt.colorbar()
cb.set_label(’log10(N)’)
plt.savefig(outfile, papertype=’A4’, format=’pdf’)
|
arxiv-papers
| 2013-09-20T14:12:04 |
2024-09-04T02:49:51.227314
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Dan Stowell and Mark D. Plumbley",
"submitter": "Dan Stowell",
"url": "https://arxiv.org/abs/1309.5275"
}
|
1309.5415
|
# Light-weight Flexible Magnetic Shields For Large-Aperture Photomultiplier
Tubes
P. DeVore D. Escontrias T. Koblesky C. J. Lin D. W. Liu K. B. Luk J.
Ngan J. C. Peng C. Polly J. Roloff H. Steiner S. Wang J. Wong M. Yeh
Department of Physics, University of California, Berkeley, CA 94720, U.S.A.
Physics Division, Lawrence Berkeley National Laboratory, Berkeley, CA 94720,
U.S.A. Department of Physics, University of Illinois at Urbana-Champaign,
Urbana, IL 61801, U.S.A. Department of Physics, Chinese University of Hong
Kong, Hong Kong, China Chemistry Department, Brookhaven National Laboratory.
Upton, NY 11973, U.S.A.
(August 20, 2013)
###### Abstract
Thin flexible sheets of high-permeability FINEMET® foils encased in thin
plastic layers have been used to shield various types of 20-cm-diameter
photomultiplier tubes from ambient magnetic fields. In the presence of the
Earth’s magnetic field this type of shielding is shown to increase the
collection efficiency of photoelectrons and can improve the uniformity of
response of these photomultiplier tubes.
###### keywords:
Magnetic Shield , Photomultiplier tube , FINEMET
###### PACS:
07.55.Nk , 85.60.Ha
, , , , , , , , , , , , ,
## 1 Introduction
It is well known that the performance of most photomultiplier tubes (PMTs) is
susceptible to even small local magnetic field such as that of the Earth,
which is typically in the range of 0.4 to 0.5 gauss [1]. The impact of this
kind of weak magnetic field can be reduced by shielding the PMT, especially
the region near the photocathode, with high-permeability materials such as mu-
metal.
We describe here a new type of magnetic shield that has been developed at the
Lawrence Berkeley National Laboratory for use in the Daya Bay Reactor Neutrino
Oscillation Experiment [2]. It utilizes FINEMET, a novel thin flexible high-
permeability material, to form a truncated cone around a large-aperture PMT
behind the photocathode. In the following sections we describe the detailed
characteristics of this shield, and present experimental results showing its
impact on the performance of some selected PMTs. In particular we will show
that it reduces variations in the PMT response caused by changes in the
relative orientation of the PMT and magnetic field, and that it significantly
improves the electron collection efficiency of the PMT. These improvements
have a direct impact on the physics objectives of an experiment in that they
result in a greater uniformity of response of the detector, and thereby yield
a more reliable measurement of the energy of an event.
## 2 Effect of Magnetic Field on PMT
We have investigated the effect of the magnetic field on the gain and charge
collection of some large-aperture PMTs. These results are of particular
relevance to experiments where the events of interest produce at most a few
photoelectrons (PEs) in a given PMT. The magnetic field also affects multi-
photoelectron events, but in this case it is not possible to separate the
effects caused by gain variation from those attributable to collection
efficiency. In this study we used Electron Tubes 9354KB [3], Hamamatsu R5912
[4], and Photonis XP1806 [5] 20-cm PMTs.
### 2.1 Effect of Magnetic Field on Collected Amount of Charge
We determined the effect of the local magnetic field on the response of a PMT
by positioning the PMT with its polar axis perpendicular to the field, as
shown in Figure 1(a). The PMT was mounted on a holder that could be rotated
about the polar axis from outside of the dark box containing the PMT. Light
from a pulsing blue LED illuminated the entire photocathode, and its intensity
could be controlled as needed in the study. The charge associated with the PMT
signal was measured with an analog-to-digital converter (ADC) as a function of
the angle between the local magnetic field and a reference on the PMT, denoted
by $\phi$. In this study, the polarized key on the socket of the PMT, shown in
Figure 1(b), was chosen as the reference direction. When this reference
direction was parallel (normal) to the local field, $\phi$ was $0^{\circ}$
($90^{\circ}$).
Figure 1: (a) Orientation of the PMT with respect to the local magnetic field
for evaluating response. (b) PMT reference direction: the polarized key
between pin 1 and pin 20 indicated by the arrow is used as the PMT reference
direction.
The $\phi$-dependence of the charge variation at some operating high voltages
for the Electron Tubes 9354KB, Hamamatsu R5912, and Photonis XP1806 PMTs is
shown in Figures 2, 3 and 4 respectively. The charge variation, $\Delta Q$, is
defined as
$\Delta Q(\phi)=\left(\frac{Q(\phi)-\overline{Q}}{\overline{Q}}\right)\times
100\%$ (1)
where $Q(\phi)$ is the amount of charge collected at angle $\phi$, and
$\overline{Q}$ is the average amount of collected charge given by
$\overline{Q}=\frac{1}{N}\sum_{i=1}^{N}Q(\phi_{i})$ (2)
with $N$ being the number of measurements done at various $\phi$ values.
For each type of PMT the result is quite insensitive to the applied voltage.
The difference in the observed behaviour between the three PMT models is a
consequence of the different designs of the dynode structure and focussing
scheme. Figure 5 shows the opening to the first dynode for these three kinds
of PMTs. Electron Tubes 9354KB has the smallest opening (most affected)
whereas Photonis XP1806 has the largest opening (least affected).
Figure 2: Effect of magnetic field on collecting charge for an unshielded
Electron Tube 9354KB PMT as a function of the angle between the PMT reference
direction and the magnetic field. Figure 3: Effect of magnetic field on
collecting charge for an unshielded Hamamatsu R5912 PMT as a function of the
angle between the PMT reference direction and the magnetic field. Figure 4:
Effect of magnetic field on collecting charge for an unshielded Photonis
XP1806 PMT as a function of the angle between the PMT reference direction and
the magnetic field. Figure 5: Opening to the first dynode. From left to right:
Electron Tubes 9354KB, Hamamatsu R5912, and Photonis XP1809.
### 2.2 Effect of Magnetic Field on Single-photoelectron (SPE) Spectrum
A revealing way to study the influence of the magnetic field on the detected
charge is to use single-photoelectron events. By making this measurement
without a magnetic shield we can directly determine how the shield affects
both the gain and the collection efficiency of the PMT. The SPE spectrum can
be obtained by reducing the light intensity on the photocathode until only
about 10% of the triggers produce a signal. From Poisson statistics about 95%
of the observed pulses are produced by SPEs. The spectrum of the charge
collected at the anode under this condition is shown in Figure 6. Since the
charge of a single electron is known a priori, the calibrated ADC channel of
the SPE peak can be used to determine the absolute gain of the PMT.
Photoelectrons generated at the photocathode are directed to the first dynode
by a set of focussing electrodes. The configuration of the focussing elements
and the effective area of the first dynode determines the collection
efficiency. The Lorentz force, v $\times$ B, on the photoelectrons in the
rather long distance between the photocathode and the first dynode causes some
of them to be swept away from the first dynode and thereby reduces the
collection efficiency. The number of events in the SPE distribution is
proportional to the efficiency of collecting the produced SPE by the first
dynode of the PMT.
Using a Hamamatsu R5912 without a magnetic shield we have obtained SPE spectra
with the PMT reference direction parallel ($\phi=0^{\circ}$) and perpendicular
($\phi=90^{\circ}$) to the local magnetic field. The observed ADC spectrum of
each orientation, a convolution of the SPE distribution and the pedestal, is
shown in Figure 6. The observed spectrum was fitted to two Gaussian functions,
one for the pedestal and the other one for the SPE distribution, and an
exponential function for modeling the non-Gaussian component of the pedestal.
The gain , peak-to-valley ratio (P/V ratio) and charge resolution of the PMT
were obtained from the best fits which are displayed in Figure 6. The results
are tabulated in Table 1. It is interesting to note that the change in gain
between the two orientations is small. From Figure 3, we see a 15% difference
in the charge variation between these two angular positions; that is, the
collection efficiency has changed by about 15%. This result suggests that the
majority of the secondary electrons produced within the dynodes are not
significantly affected by the magnetic field. However, the peak-to-valley
ratio gets worse when the magnetic effect is more pronounced. A series of
measurements was made to determine how the collection efficiency in the
presence of a magnetic field depends on the potential difference between the
photocathode and first dynode. The results showed almost no effect on the
collection efficiency when the potential difference was within 250 V of its
nominal value, typically several hundred volts, for a gain of $10^{7}$.
Figure 6: Magnetic effect on the SPE distribution of a Hamamatsu R5912 without shielding. Figure on the left is obtained with the PMT reference direction parallel to the magnetic field ($\phi=0^{\circ}$); figure on the right is determined with the magnetic field perpendicular to the PMT reference direction ($\phi=90^{\circ}$). The red curve in each figure is the best fit to the spectrum. Table 1: SPE Spectra of a Hamamatsu R5912 in the presence of magnetic field. PMT Orientation | Gain | P/V Ratio | Charge Resolution
---|---|---|---
$\varparallel$ magnetic field ($\phi=0^{\circ}$) | $(1.030\pm 0.003)\times 10^{7}$ | 2.77 $\pm$ 0.12 | 0.343 $\pm$ 0.003
$\perp$ magnetic field ($\phi=90^{\circ}$) | $(1.043\pm 0.004)\times 10^{7}$ | 2.13 $\pm$ 0.11 | 0.333 $\pm$ 0.003
## 3 FINEMET magnetic shield
In this section we describe the magnetic shield made of FINEMET, the material
we used for reducing the effect of magnetic field on the Electron Tubes
9354KB, Hamamatsu R5912, and Photonis XP1806 PMTs.
### 3.1 FINEMET®
FINEMET is a trademark of Hitachi Materials Ltd. for an amorphous metallic
alloy ribbon made up of nano-crystalline grains of iron, silicon, boron, and
small amounts of copper and niobium [6]. As shown in Table 2, this novel soft
magnetic alloy has high saturation magnetic flux density, low core loss, and
high permeability over a wide range of frequencies. Its permeability varies by
less than $\pm$10% between $-50~{}^{\circ}$C and $150~{}^{\circ}$C, with very
small aging effects. We used 0.12-mm-thick FINEMET sheets for fabricating the
magnetic shields. This thin flexible material is made of a layer of 18
$\mu$m-thick FINEMET FT-3M tape laminated with 25-$\mu$m-thick hot-melt
adhesive between two polyethylene terephthalate (PET) films. It can operate in
a temperature range between $-40~{}^{\circ}$C and $80~{}^{\circ}$C. It is
interesting to note that one layer of FINEMET can reduce the local magnetic
field of $4\times 10^{-5}$ T by about 20 dB (from Figure 4 of reference [7]).
Property | Value
---|---
Density | 7300 kg/m3
Resistivity | 1.2 $\mu\Omega$-m
Curie temperature | 570 ∘C
Magnetic flux density (DC H = 800 A/m) | 1.23 T
Maximum relative permeability (DC) | 70,000
Table 2: Selected physical properties of Hitachi FINEMET FT-3M [7].
In addition to the physical properties provided by the manufacturer and listed
in Table 2, we have measured the amount of natural radioactivity of a few
samples of FINEMET. Our results were $(18\pm 5)$ ppb for the 238U series,
$(15\pm 5)$ ppb for the 232Th series, and $(4\pm 2)$ ppm for K, which
translate to decay rates of about 0.2 Bq/kg, 0.06 Bq/kg, and 0.1 Bq/kg for
238U, 232Th, and 40K respectively. We have observed that if it is not stored
in a ventilated location progenies of radon can accumulate gradually on the
surfaces of the laminated FINEMET foil.
Besides radioactivity, we have also investigated the compatibility of FINEMET
with mineral oil and 18-M$\Omega$ ultra-pure water, which are commonly used in
nuclear and particle-physics experiments. By submerging a piece of FINEMET
sample with known dimensions in each liquid at elevated temperature for
accelerated aging studies, the absorbance of the liquid was monitored as a
function of time. For mineral oil at 40 ∘C, with a ratio of contact surface to
liquid volume (S/V) of 0.39 cm-1, its absorbance as a function of wavelength
and time is shown in Figure 7. The change in absorbance due to the presence of
FINEMET is minute for wavelengths between 300 nm and 650 nm. The absorbance of
purified water at 70 ∘C for a FINEMET sample with a S/V of 0.32 cm-1 is
summarized in Figure 8. In the vicinity of 350 nm, the rate of change in
absorbance levels off in about two weeks after a relatively fast degradation
in the beginning. Note that the integrity of FINEMET has not changed in both
liquids.
Figure 7: Absorbance of mineral oil at $40~{}^{\circ}$C as a function of
wavelength for a submerged FINEMET sample. Figure 8: Absorbance of purified
water at $70~{}^{\circ}$C as a function of wavelength for a submerged FINEMET
sample.
### 3.2 Fabrication of Magnetic Shield
As illustrated in Figure 9, a truncated conical magnetic shield for the
Hamamatsu R5912 PMTs used in the Daya Bay Experiment was realized by cutting a
two-dimensional fan out of a piece of FINEMET foil with a water jet. The
conically shaped shield was made by wrapping the fan around the PMT as shown
in Figure 10. It was firmly held in place between the equator of the glass
bulb and the acrylic potting shell surrounding the base by tucking the tabs
into the pre-cut slots. The shield is flexible, so any deformation of the cone
can be taken out easily by simple manipulation of the surface. The weight of
such a FINEMET shield for a typical 20-cm PMT is about 20 g.
Figure 9: Two-dimensional pattern of an unfolded truncated conical FINEMET
magnetic shield. Figure 10: Cross section of a Hamamatsu R5912 PMT used in
Daya Bay with a truncated conical magnetic shield. The dashed line is the
photocathode.
### 3.3 Performance of the FINEMET magnetic shield
The effectiveness of the truncated conical FINEMET shield depends on the slant
height of the cone. For the Hamamatsu R5912 assembly used in the Daya Bay
experiment, the magnetic shield extends from the acrylic sealing shell to the
equator of the bulb, where the photocathode ends. We have evaluated the
performance of the FINEMET magnetic shields with slant heights of 15.4 cm,
16.7 cm, 18.1 cm, 19.2 cm, and 23.0 cm, where the 23 cm-long shield extends to
the plane tangent to the apex of the PMT bulb. The nominal slant height of the
shield used in Daya Bay is 15.4 cm, ending right at the contact point of the
bulb near the equator. For the Photonis XP1806 PMT assembly we used shields
with slant heights of 14.3 cm and 23.6 cm which extended to the equator and
the apex of the PMT bulb respectively. The corresponding slant heights are
14.5 cm and 22.5 cm for the Electron Tubes 5394KB PMT.
#### 3.3.1 Improvement on the collection efficiency
With the polar axis of the PMT normal to the local magnetic field as shown in
Figure 1, the amount of collected charge relative to the one without shielding
at $\phi=0^{\circ}$ (relative collection efficiency) for the Hamamatsu R5912
PMT as a function of the slant height of the FINEMET shield for different
$\phi$ angles is shown in Figure 11. The magnetic shield clearly improves the
collection of charge. Relative to the case of without any shielding, a FINEMET
shield with a slant height of 15.4 cm (23.0 cm) improves the average
collection efficiency by 8% (19%). Furthermore, as shown in Figure 12, the
charge variation as a function of the $\phi$ angle is reduced when the slant
height increases. For example, without shielding, the peak-to-peak difference
in the charge variation is about 23%. A FINEMET shield with a slant height of
23.0 cm shrinks this difference to approximately 5%. Even a shield that covers
only up to the equator (15.4 cm) allows us to achieve almost the same
performance. This implies that the shield provides a less non-uniform
collection efficiency to this level over the entire photocathode.
Similarly, the results on collecting charge and minimizing the variation of
the collection efficiency for the Photonis XP1806 and Electron Tubes 9354KB
PMTs with FINEMET shields are depicted in Figures 13, 14, 15 and 16. For the
Photonis XP1806 PMT, on average, the collection efficiency is increased by
about 5% (14%) for a shield with a slant height of 14.3 cm (23.6 cm), and the
variation in the amount of charge collected is reduced from 10% to 1.3% which
is close to our systematic limit. The most dramatic improvement is achieved
with the Electron Tubes 9354KB PMT. In this case the average collection
efficiency is up by about 25% for a FINEMET shield extending to the equator
and 40% for shields reaching to the apex of the PMT bulb. The variation in the
collection efficiency is reduced from 46% with no shield to 21% (5%) for a
shield with a slant height of 14.5 cm (22.5 cm).
Figure 11: Relative collection efficiency as a function of the slant height of
truncated conical FINEMET shield for a Hamamatsu R5912 PMT oriented at
different $\phi$ angles. Figure 12: Charge variation due to magnetic field for
Hamamatsu R5912 PMT with FINEMET shields of different slant heights. Figure
13: Relative collection efficiency as a function of slant height of truncated
conical FINEMET shield for a Photonis XP1806 PMT oriented at different $\phi$
angles. Figure 14: Charge variation due to magnetic field for Photonis XP1806
PMT with FINEMET shields of different slant heights. Figure 15: Relative
collection efficiency as a function of slant height of truncated conical
FINEMET shield for an Electron Tubes 9354KB PMT oriented at different $\phi$
angles. Figure 16: Charge variation due to magnetic field for Electron Tubes
9354KB PMT with FINEMET shields of different slant heights.
#### 3.3.2 Improvement on the SPE Spectrum
In Table 3 we summarize the results on the SPE spectrum for the Hamamatsu
R5912 PMT with FINEMET shields of different sizes. Compared to the results for
the case without shielding (Table 1), we see that the shield can improve the
gain by 5%, P/V ratio by up to 17% (34%) for $\phi=0^{\circ}$ ($90^{\circ}$),
and charge resolution by less than 10%. Furthermore, the variation of the P/V
ratio due to different orientation of the PMT relative to the magnetic field
is significantly reduced. Consequently the systematic uncertainty in the
charge measurement is minimized.
Table 3: Characteristics of the SPE Spectrum for a Hamamatsu R5912 PMT with FINEMET magnetic shields of different slant heights. PMT Reference Direction | Gain | P/V Ratio | Charge Resolution
---|---|---|---
Slant height = 15.4 cm
$\varparallel$ magnetic field ($\phi=0^{\circ}$) | $(1.080\pm 0.003)\times 10^{7}$ | 3.27 $\pm$ 0.18 | 0.316 $\pm$ 0.002
$\perp$ magnetic field ($\phi=90^{\circ}$) | $(1.092\pm 0.003)\times 10^{7}$ | 2.86 $\pm$ 0.14 | 0.324 $\pm$ 0.003
Slant height = 18.1 cm
$\varparallel$ magnetic field ($\phi=0^{\circ}$) | $(1.091\pm 0.003)\times 10^{7}$ | 3.17 $\pm$ 0.16 | 0.311 $\pm$ 0.002
$\perp$ magnetic field ($\phi=90^{\circ}$) | $(1.088\pm 0.03)\times 10^{7}$ | 3.03 $\pm$ 0.15 | 0.318 $\pm$ 0.003
Slant height = 23.0 cm
$\varparallel$ magnetic field ($\phi=0^{\circ}$) | $(1.081\pm 0.003)\times 10^{7}$ | 3.20 $\pm$ 0.16 | 0.315 $\pm$ 0.002
$\perp$ magnetic field ($\phi=90^{\circ}$) | $(1.085\pm 0.003)\times 10^{7}$ | 2.93 $\pm$ 0.14 | 0.311 $\pm$ 0.002
#### 3.3.3 Position-dependence of PMT Response
We determined the response of a Hammamatsu R5912 PMT at each point of a grid
on the photocathode by measuring the PE spectrum with a pulsing LED [8]. The
orientation of the PMT is shown in Figure 17. We split the light pulse from
the LED into two parts with optical fibers. One part of the LED light was sent
to a reference PMT to monitor the intensity of the LED, and the other part to
the photocathode of the Hamamatsu PMT. A collimator was used to limit the spot
size of the incident light at the photocathode to a circle with a radius of 2
mm. The LED intensity was set to produce at least 20 photoelectrons per pulse
on the photocathode. Runs of ten seconds duration at a pulse rate of 100 Hz
were used to measure the response of the PMT for every point on the
photocathode. An ADC was gated by the pulse generator used to drive the LED. A
computer recorded the integrated counts in each ADC channel for each run. The
light source was moved under computer control to the next point on the grid
until the whole surface of the photocathode was scanned.
Figure 17: (a) Orientation of the photocathode of a Hamamatsu R5912 PMT with
respect to the local magnetic field. The x-y plane is normal to the polar axis
of the PMT. The magnetic field is along the negative-y direction. (b)
Position-dependent response of a Hamamatsu R5912 PMT with no magnetic shield.
The colored scale is the average number of detected PEs. Figure 18: Position-
dependent response of a Hamamatsu R5912 PMT with a FINEMET shield (a) with a
slant height of 15.4 cm, (b) with a slant height of 20.5 cm. The colored scale
is the average number of detected PEs.
In the absence of any magnetic field we expect rather uniform collection
efficiency over the photocathode. However, the Earth’s magnetic field in the
transverse direction creates a position-dependent response. In this case the
field acting on a PMT with no shielding produces large non-uniformity as shown
in Figure 17. The lowest collection efficiency occurs in a region centered at
x $\approx$ 6 cm and y = 0 cm. In addition, the response of the PMT is almost
symmetric about the x-axis. This result is a direct consequence of the fact
that the internal structure of the PMT is symmetric about the x-axis, so that
with the magnetic field in the negative y-direction the response of the PMT
should be symmetric across the x-axis. This symmetry should not exist for the
y-axis because the internal structure of the PMT is different across the
y-axis, and also because a magnetic field in the negative y-direction will
affect the motion of electrons in the x-direction.
A FINEMET shield around the PMT makes the response more uniform over the whole
photocathode. The results for the shields with slant height of 15.4 cm and
20.5 cm are presented in Figure 18. The 20.5-cm shield has a more uniform
response than the 15.4-cm shield, although even with this long shield there
are still significant residual non-uniformities. However, both of these
shields clearly produce a better performance for the PMT than that obtained
with no shield.
## 4 Conclusions
We have presented the performance of a truncated conical magnetic shield made
of FINEMET. Mechanically, it is thin, light, flexible, relative cheap, and
easy to handle. In addition, the laminated FINEMET sheet is compatible with
mineral oil and purified water. We have demonstrated that FINEMET is very
effective in shielding large-aperture PMTs against magnetic field. As a
result, the collection efficiency of the PMT is improved and the response of
the PMT is more uniform across the entire photocathode.
## 5 Acknowledgements
We would like to express our gratitude to the technical staffs of the Lawrence
Berkeley National Laboratory for their excellent support. This work was
partially supported by the Director, Office of Science, Office of Basic Energy
Sciences, Office of High Energy Physics, of the U.S. Department of Energy
under Contract No. DE-AC02-05CH11231and DE-AC02-98CH10886, the U.S. National
Science Foundation, the Research Grant Council of the Hong Kong Special
Administrative Region, China (Project Nos. CUHK 1/07C and CUHK3/CRF/10). J.N.
and J. W. were also partially supported by the Department of Physics, Chinese
University of Hong Kong.
## References
* [1] Hamamatsu Photonics, ”Photomultiplier Tubes: Basics and Applications”, 3rd Edition, (2007).
* [2] X. Guo et. al. (Daya Bay Collaboration), hep-ex/0701029 (2007).
* [3] Electron Tubes 9354KB Data Sheet, Electron Tubes Inc., Sept., 2002.
* [4] Hamamatsu R5912 Data Sheet, Hamamatsu Co., Sept., 1998.
* [5] Photonis Hemispherical PMT Data Sheet, Photonis Co, Sept., 2006.
* [6] R. Lebourgeois et al., J. Magn. Magn. Mater. 254-255, 191 (2003).
* [7] Catalog number HL-FM16-F, Hitachi Metals, Ltd. (http://www.hitachi-metals.co.jp).
* [8] T. Koblesky et. al., Nucl. Instrum. Method A670, 40 (2012).
|
arxiv-papers
| 2013-09-20T23:45:35 |
2024-09-04T02:49:51.241812
|
{
"license": "Public Domain",
"authors": "P. DeVore, D. Escontrias, T. Koblesky, C. J. Lin, D. W. Liu, K. B.\n Luk, J. Ngan, J. C. Peng, C. Polly, J. Roloff, H. Steiner, S. Wang, J. Wong\n and M. Yeh",
"submitter": "Kam Biu Luk",
"url": "https://arxiv.org/abs/1309.5415"
}
|
1309.5427
|
# Latent Fisher Discriminant Analysis
Gang Chen
Department of Computer Science and Engineering
SUNY at Buffalo
[email protected]
###### Abstract
Linear Discriminant Analysis (LDA) is a well-known method for dimensionality
reduction and classification. Previous studies have also extended the binary-
class case into multi-classes. However, many applications, such as object
detection and keyframe extraction cannot provide consistent instance-label
pairs, while LDA requires labels on instance level for training. Thus it
cannot be directly applied for semi-supervised classification problem. In this
paper, we overcome this limitation and propose a latent variable Fisher
discriminant analysis model. We relax the instance-level labeling into bag-
level, is a kind of semi-supervised (video-level labels of event type are
required for semantic frame extraction) and incorporates a data-driven prior
over the latent variables. Hence, our method combines the latent variable
inference and dimension reduction in an unified bayesian framework. We test
our method on MUSK and Corel data sets and yield competitive results compared
to the baseline approach. We also demonstrate its capacity on the challenging
TRECVID MED11 dataset for semantic keyframe extraction and conduct a human-
factors ranking-based experimental evaluation, which clearly demonstrates our
proposed method consistently extracts more semantically meaningful keyframes
than challenging baselines.
## 1 Introduction
Linear Discriminant Analysis (LDA) is a powerful tool for dimensionality
reduction and classification that projects high-dimensional data into a low-
dimensional space where the data achieves maximum class separability [10, 12].
The basic idea in classical LDA, known as the Fisher Linear Discriminant
Analysis (FDA) is to obtain the projection matrix by minimizing the within-
class distance and maximizing the between-class distance simultaneously to
yield the maximum class discrimination. It has been proved analytically that
the optimal transformation is readily computed by solving a generalized
eigenvalue problem [12]. In order to deal with multi-class scenarios [10], LDA
can be easily extended from binary-case and generally used to find a subspace
with $d-1$ dimensions for multi-class problems, where $d$ is the number of
classes in the training dataset. Because of its effectiveness and
computational efficiency, it has been applied successfully in many
applications, such as face recognition [3] and microarray gene expression data
analysis. Moreover, LDA was shown to compare favorably with other supervised
dimensionality reduction methods through experiments [20].
However, LDA expects instance/label pairs which are surprisingly prohibitive
especially for large training data. In the last decades, semi-supervised
methods have been proposed to utilize unlabeled data to aid classification or
regression tasks under situations with limited labeled data, such as
Transductive SVM (TSVM) [21, 14] and Co-Training [4]. Correspondently, it is
reasonable to extend the supervised LDA into a semi-supervised method, and
many approaches [5, 26, 20] have been proposed. Most of these methods are
based on transductive learning. In other words, they still need instance/label
pairs. However, many real applications require bag-level labeling [1], such as
object detection [11] and event detection [19].
In this paper, we propose a Latent Fisher Discriminant Analysis model (or LFDA
in short), which generalizes Fisher LDA model [10]. Our model is inspired by
MI-SVM [1] or latent SVM [11] and multiple instance learning problems [8, 17].
On the one hand, recently applications in image and video analysis require a
kind of bag-level label. Moreover, using latent variable model for this kind
of problem shows great improvement on object detection [11]. On the other
hand, the requirement of instance/label pairs in the training data is
surprisingly prohibitive especially for large training data. The bag-level
labeling methods are a good solution to this problem.
MI-SVM or latent SVM is a kind of discriminative model by maximizing a
posterior probability. Our model unify the discriminative nature of the Fisher
linear discriminant with a data driven Gaussian mixture prior over the
training data in the Bayesian framework. By combining these two terms in one
model, we infer latent variables and projection matrix in alternative way
until convergence. We demonstrate this capability on MUSK and Corel data sets
for classification, and on TRECVID MED11 dataset for keyframe extraction on
five video events [19].
## 2 Related Work
Linear Discriminant Analysis (LDA) has been a popular method for dimension
reduction and classification. It searches a projection matrix that
simultaneously maximizes the between-class dissimilarity and minimizes the
within-class dissimilarity to increase class separability, typically for
classification applications. LDA has attracted an increasing amount of
attention in many applications because of its effectiveness and computational
efficiency. Belhumeur et al proposed PCA+LDA [3] for face recognition. Chen et
al projects the data to the null space of the within-class scatter matrix and
then maximizes the between-class scatter in this space [6] to deal with the
situation when the size of training data is smaller than the dimensionality of
feature space. [22] combines the ideas above, maximizes the between-class
scatter matrix in the range space and the null space of the within-class
scatter matrix separately and then integrates the two parts together to get
the final transformation. [25] is also a two-stage method which can be divided
into two steps: first project the data to the range space of the between-class
scatter matrix and then apply LDA in this space. To deal with non-linear
scenarios, the kernel approach [21] can be applied easily via the so-called
kernel trick to extend LDA to its kernel version, called kernel discriminant
analysis [2], that can project the data points nonlinearly. Recently, sparsity
induced LDA is also proposed [18].
However, many real-world applications only provide labels on bag-level, such
as object detection and event detection. LDA, as a classical supervised
learning method, requires a training dataset consisting of instance and label
pairs, to construct a classifier that can predict outputs/labels for novel
inputs. However, directly casting LDA as a semi-supervised method is
challenging for multi-class problems. Thus, in the last decades, semi-
supervised methods become a a hot topic. One of the main trend is to extend
the supervised LDA into a semi-supervised method [5, 26, 20], which attempts
to utilize unlabeled data to aid classification or regression tasks under
situations with limited labeled data. [5] propose a novel method, called Semi-
supervised Discriminant Analysis, which makes use of both labeled and
unlabeled samples. The labeled data points are used to maximize the
separability between different classes and the unlabeled data points are used
to estimate the intrinsic geometric structure of the data. [20] propose a
semi-supervised dimensionality reduction method which preserves the global
structure of unlabeled samples in addition to separating labeled samples in
different classes from each other. Most of these semi-supervised methods model
the geometric relationships between all data points in the form of a graph and
then propagate the label information from the labeled data points through the
graph to the unlabeled data points. Another trend prefers to extent LDA into
an unsupervised senarios. For example, Ding et al propose to combine LDA and
K-means clustering into the LDA-Km algorithm [9] for adaptive dimension
reduction. In this algorithm, K-means clustering is used to generate class
labels and LDA is utilized to perform subspace selection.
Our solution is a new latent variable model called Latent Fisher Discriminant
Analysis (LFDA), which complements existing latent variable models that have
been popular in the recent vision literature [11] by making it possible to
include the latent variables into Fisher discriminant analysis model. Unlike
previous latent SVM [11] or MI-SVM [1] model, we extend it with prior data
distribution to maximize a joint probability when inferring latent variables.
Hence, our method combines the latent variable inference and dimension
reduction in an unified Bayesian framework.
## 3 Latent Fisher discriminant analysis
We propose a LFDA model by including latent variables into the Fisher
discriminant analysis model. Let ${\mathcal{X}}=\\{{\bf x_{1}},{\bf
x_{2}},...,{\bf x_{n}}\\}$ represent $n$ bags, and the corresponding labels
${\mathcal{L}}=\\{l_{1},l_{2},...,l_{n}\\}$. For each ${\bf
x_{i}}\in\mathcal{X}$, ${\bf x_{i}}$ can be treated as a bag (or video) in
[1], and its label $l_{i}$ is categorical and assumes values in a finite set,
e.g. $\\{1,2,...,C\\}$. Let ${\bf x_{i}}\in\mathbb{R}^{d\times n_{i}}$, which
means it contains $n_{i}$ instances (or frames), ${\bf
x_{i}}=\\{x_{i}^{1},x_{i}^{2},...,x_{i}^{n_{i}}\\}$, and its $j^{th}$ instance
is a vector in $\mathbb{R}^{d}$, namely $x_{i}^{j}\in\mathbb{R}^{d}$. Fisher’s
linear discriminant analysis pursue a subspace $\mathcal{Y}$ to separate two
or more classes. In other words, for any instance $x\in\mathcal{X}$, it
searches for a projection $f:x\rightarrow y$, where
$y\in\mathbb{R}^{d^{\prime}}$ and $d^{\prime}\leq d$. In general, $d^{\prime}$
is decided by $C$, namely $d^{\prime}=C-1$. Suppose the projection matrix is
$\mathcal{P}$, and $y=f(x)=\mathcal{P}x$, then latent Fisher LDA proposes to
minimize the following ratio:
$\displaystyle\centering({\mathcal{P}^{*}})=\operatornamewithlimits{arg\,min}_{\mathcal{P},z}J({\mathcal{P},z})=\operatornamewithlimits{arg\,min}_{\mathcal{P},z}\textrm{trace}\left(\frac{\mathcal{P}^{T}\Sigma_{w}(x,z)\mathcal{P}}{\mathcal{P}^{T}\Sigma_{b}(x,z)\mathcal{P}}+\beta\mathcal{P}^{T}\mathcal{P}\right)\@add@centering$
(1)
where $z$ is the latent variable, $\beta$ is an weighting parameter for
regularization term. The set $z\in Z(x)$ defines the possible latent values
for a sample $x$. In our case, $z\in\\{1,2,...,C\\}$. $\Sigma_{b}(x,z)$ is
between class scatter matrix and $\Sigma_{w}(x,z)$ is within class scatter
matrix. However, LDA is dependent on a categorical variable (i.e. the class
label) for each instance $x$ to compute $\Sigma_{b}$ and $\Sigma_{w}$. In our
case, we only know bag-level labels, not on instance-level labels. To minimize
$J({\mathcal{P}},z)$, we need to solve $z(x)$ for any given $x$. This problem
is a chicken and egg problem, and can be solved by alternating algorithms,
such as EM [7]. In other words, solve $\mathcal{P}$ in Eq. (1) with fixed $z$,
and vice versa in an alternating strategy.
### 3.1 Updating $z$
Suppose we have found the projection matrix $\mathcal{P}$, and corresponding
subspace $\mathcal{Y}=\mathcal{P}\mathcal{X}$, where $\mathcal{Y}=\\{{\bf
y_{1}},{\bf y_{2}},...,{\bf y_{n}}\\}$. Instead of inferring latent variables
at instance-level in latent SVM, we propose latent variable inference at
clustering-level in the projected space $\mathcal{Y}$. That means all elements
in the same cluster have same label. Such assumption is reasonable because
elements in the same cluster are close to each other. On the other hand,
cluster-level inference can speed up the learning process. We extend mixture
discriminative analysis model in [13] by incorporating latent variables over
all instances for an given class. As in [13], we assume each class $i$ is a
$K$ components of Gaussians,
$\centering
p(x|\lambda_{i})=\sum_{j=1}^{K}\pi_{i}^{j}g(x|\mu_{i}^{j},\Sigma_{i}^{j})\@add@centering$
(2)
where $x$ is a $d$-dimensional continuous-valued data vector (i.e. measurement
or features); $\pi_{i}=\\{\pi_{i}^{j}\\}_{j=1}^{K}$ are the mixture weights,
and $g(x|\mu_{i}^{j},\Sigma_{i}^{j})$, $j\in[1,K]$, are the component Gaussian
densities with $\mu_{i}=\\{\mu_{i}^{j}\\}_{j=1}^{K}$ as mean and
$\Sigma_{i}=\\{\Sigma_{i}^{j}\\}_{j=1}^{K}$ as covariance. $\lambda_{i}$ is
the parameters for class $i$ which we need to estimate,
$\lambda_{i}=\\{\pi_{i},\mu_{i},\Sigma_{i}\\}$.
Hence, for each class $i\in\\{1,2,\dots,C\\}$, we can estimate $\lambda_{i}$
and get its $K$ subsets $S_{i}=\\{S_{i}^{1},S_{i}^{2},...,S_{i}^{K}\\}$ with
EM algorithm. Suppose we have the discriminative weights (or posterior
probability) for the $K$ centers in each class,
$w_{i}=\\{w_{i}^{1},w_{i}^{2},...,w_{i}^{K}\\}$, which are the posterior
probability determined by the latent FDA and will be discussed later. We
maximize one of the following two equations:
Maximizing a posterior probability:
$\displaystyle\centering\mu_{i}^{j}=\operatornamewithlimits{arg\,max}_{\mu_{i}^{j}\in\mu_{i},j\in[1,K]}w_{i}=\operatornamewithlimits{arg\,max}_{\mu_{i}^{j}\in\mu_{i},j\in[1,K]}p(z_{i}|\mu_{i},\mathcal{P})\@add@centering$
(3a)
Maximizing the joint probability with prior:
$\displaystyle\centering\mu_{i}^{j}=\operatornamewithlimits{arg\,max}_{\mu_{i}^{j}\in\mu_{i},j\in[1,K]}(\pi_{i}\circ
w_{i})\@add@centering$ (3b)
where $z_{i}$ is the latent label assignment, $\pi_{i}$ is the prior
clustering distributions for $\lambda_{i}$ in class $i$, $w_{i}$ is the
posterior (or weight) determined by kNN voting (see further) in the subspace
and $\circ$ is the pointwise production or Hadamard product. We treat Eq. (3a)
as the latent Fisher discriminant analysis model (LFDA), because it takes the
same strategy as the latent SVM model [1, 11]. As for Eq. (3b), we extend LFDA
by combining the both factors (representative and discriminative) together,
and find the cluster $S_{i}^{j}$ in class $i$ by maximizing Eq. (3b). In a
sense, Eq. (3b) considers the prior distribution from the training dataset,
thus, we treat it as the joint latent Fisher discriminant analysis model
(JLFDA) or LFDA with prior. In the nutshell, we propose a way to formulate
discriminative and generative methods under the unified Bayesian framework. We
comparatively analyze both of these models (Section 4).
Consequently, if we select the cluster $S_{i}^{j}$ with the mean $\mu_{i}^{j}$
which maximizes the above equation for class $i$, we can relabel all samples
$x$ positive for class $i$ and the rest negative, subject to
$y={\mathcal{P}}x$ and $y\in S_{i}^{j}$. Then, we construct a new training
data ${\mathcal{X^{+}}}=\\{{\bf x_{1}^{+}},{\bf x_{2}^{+}},...,{\bf
x_{n}^{+}}\\}$, with labels ${\mathcal{L^{+}}}=\\{{\bf z_{1}^{+}},{\bf
z_{2}^{+}},...,{\bf z_{n}^{+}}\\}$, where ${\bf x_{i}^{+}}=S_{i}^{j}$ for
class $i$ with $n_{i}^{j^{\prime}}$ elements, and its labels ${\bf
z_{i}^{+}}=\\{z_{i}^{1},z_{i}^{2},...,z_{i}^{n_{i}^{j^{\prime}}}\\}$ on
instance level. Obviously, ${\bf x_{i}^{+}}\subseteq{\bf x_{i}}$ and
${\mathcal{X^{+}}}$ is a subset of ${\mathcal{X}}$. The difference between
${\mathcal{X^{+}}}$ and ${\mathcal{X}}$ lies that every element
$x_{i}^{+}\in{\bf x_{i}^{+}}$ has label $z(x_{i}^{+})$ decided by Eq.
(LABEL:eq:Eq3sub), while ${\bf x_{i}}\subset\mathcal{X}$ only has bag level
label.
### 3.2 Updating projection $\mathcal{P}$
When we have labels for the new training data ${\mathcal{X^{+}}}$, we use the
Fisher LDA to minimize $J({\mathcal{P}},z)$. Note that Eq. (1) is invariant to
the scale of the vector $\mathcal{P}$. Hence, we can always choose
$\mathcal{P}$ such that the denominator is simply
$\mathcal{P}^{T}\Sigma_{b}\mathcal{P}={\bf 1}$. For this reason we can
transform the problem of minimizing Eq. (1) into the following constrained
optimization problem [10, 12, 24]:
$\displaystyle\centering{\mathcal{P}^{*}}\@add@centering$
$\displaystyle=\operatornamewithlimits{arg\,min}_{{\mathcal{P}}}\textrm{trace}\left(\mathcal{P}^{T}\Sigma_{w}(x,z)\mathcal{P}+\beta\mathcal{P}^{T}\mathcal{P}\right)$
$\displaystyle\quad\textrm{s. t.
}{\mathcal{P}^{T}\Sigma_{b}(x,z)\mathcal{P}}={\bf 1}$ (4)
where 1 is the identity matrix in $\mathbb{R}^{d^{\prime}\times d^{\prime}}$.
The optimal Multi-class LDA consists of the top eigenvectors of
$(\Sigma_{w}(x,z)+\beta)^{\dagger}\Sigma_{b}(x,z)$ corresponding to the
nonzero eigenvalues [12], here $(\Sigma_{w}(x,z)+\beta)^{\dagger}$ denotes the
pseudo-inverse of $\Sigma_{w}(x,z)+\beta$. After we calculated $\mathcal{P}$,
we can project ${\mathcal{X^{+}}}$ into subspace ${\mathcal{Y^{+}}}$. Note
that in the subspace ${\mathcal{Y^{+}}}$, any $y^{+}\in{\mathcal{Y^{+}}}$
preserves the same labels as in the original space. In other words,
${\mathcal{Y^{+}}}$ has corresponding labels ${\mathcal{L^{+}}}$ at element
level, namely $z(y^{+})=z(x^{+})$.
In general, multi-class LDA [24] uses kNN to classify new input data. We
compute $w_{i}$ using the following kNN strategy: for each sample
$x\in{\mathcal{X}}$, we get $y=\mathcal{P}x$ by projecting it into subspace
${\mathcal{Y}}$. Then, for $y\in{\mathcal{Y}}$, we choose its $N$ nearest
neighbors from ${\mathcal{Y^{+}}}$, and use their labels to voting each
cluster $S_{i}^{j}$ in each class $i$. Then, we compute the following
posterior probability:
$\displaystyle\centering w_{i}^{j}\@add@centering$
$\displaystyle=p(z_{i}=1|\mu_{i}^{j})=p(\mu_{i}^{j}|z_{i}=1)p(z_{i}=1)$
$\displaystyle=p(z_{i}=1)\frac{p(\mu_{i}^{j},z_{i}=1)}{\sum_{i=1}^{C}p(\mu_{i}^{j},z_{i}=1)}$
(5)
It counts all $y\in{\mathcal{Y}}$ fall into $N$ nearest neighbor of
$\mu_{i}^{j}$ with label $z_{i}$. Note that kNN is widely used as the
classifier in the subspace after LDA transformation. Thus, Eq. (3.2) consider
all training data to vote the weight for each discriminative cluster
$S_{i}^{j}$ in every class $i$. Hence, we can find the most discriminative
cluster $S_{i}^{j}$, s.t. $w_{i}^{j}>w_{i}^{k}$, $k\in[1,K],k\neq j$.
Algorithm. We summarize the above discussion in pseudo code. To put simply, we
update ${\mathcal{P}}$ and $z$ in an alternative manner, and accept the new
projection matrix ${\mathcal{P}}$ with LDA on the relabeled instances. Such
algorithm can always convenge in around 10 iterations. After we learned matrix
$\mathcal{P}$ and $\\{\lambda_{i}\\}_{i=1}^{C}$ by maximizing Eq. (3), we can
use them to select representative and discriminative frames from video
datasets by nearest neighbor searching.
Algorithm 1
Input: training data ${\mathcal{X}}$ and its labels $\mathcal{L}$ at video
level, $\beta$, $K$, $N$, $T$ and $\epsilon$.
Output: $\mathcal{P}$, $\\{\lambda_{i}\\}_{i=1}^{C}$
1: Initialize $\mathcal{P}$ and $w_{i}$;
2: for $Iter=1$ to $T$ do
3: for $i=1;i<=C;i++$ do
4: Project all the training data ${\mathcal{X}}$ into subspace ${\mathcal{Y}}$
using ${\mathcal{Y}}$=$\mathcal{P}{\mathcal{X}}$;
5: For each class $i\in[1,C]$, using Gaussian mixture model to partition its
elements in the subspace, and compute $\lambda_{i}$ =
$\\{\pi_{i},\mu_{i},\Sigma_{i}\\}$;
6: Maximize Eq. (3) to find $S_{i}^{j}$ with center $\mu_{i}^{j}$;
7: Relabel all elements positive in the cluster $S_{i}^{j}$ for class $i$
according to Eq. (LABEL:eq:Eq3sub);
8: end for
9: Update $z$ and construct the new subset ${\mathcal{X^{+}}}$ and its labels
$\mathcal{L}^{+}$ for all $C$ classes;
10: Do Fisher linear discriminant analysis and update $\mathcal{P}$
11: if $\mathcal{P}$ converge (change less than $\epsilon$), then break
12: Compute $N$ nearest neighbors for each training data, and calculate
discriminative weight $w_{i}$ for each class $i$ according to Eq. (3.2).
13: end for
14: Return ${\mathcal{P}}$ and cluster centers $\\{\lambda_{i}\\}_{i=1}^{C}$
learned respectively for all $C$ classes;
### 3.3 Convergence analysis
Our method updates latent variable $z$ and then $\mathcal{P}$ in an
alternative manner. Such strategy can be attributed to the hard assignment of
EM algorithm. Recall the EM approach:
$\displaystyle\centering\mathcal{P^{*}}=\operatornamewithlimits{arg\,max}_{\mathcal{P}}p(\mathcal{X},\mathcal{L}|\mathcal{P})=\operatornamewithlimits{arg\,max}_{\mathcal{P}}\sum_{i=1}^{C}p(\mathcal{X},\mathcal{L}|z_{i})p(z_{i}|\mathcal{P})\@add@centering$
(6)
then the likelihood can be optimized using iterative use of the EM algorithm.
###### Theorem 3.1
Assume the latent variable $z$ is inferred for each instance in $\mathcal{X}$,
then to maximize the above function is equal to maximize the following
auxiliary function
$\displaystyle\centering\mathcal{P}=\operatornamewithlimits{arg\,max}_{\mathcal{P}}\sum_{i=1}^{C}p(z_{i}|\mathcal{X},\mathcal{L},\mathcal{P}^{\prime})ln\bigg{(}p(\mathcal{X},\mathcal{L}|z_{i})p(z_{i}|\mathcal{P})\bigg{)}\@add@centering$
(7)
This proof can be shown using Jensen’s inequality.
###### Lemma 3.2
The hard assignment of latent variable $z$ by maximizing Eq. (3) is a special
case of EM algorithm.
* Proof
$\displaystyle\sum_{i=1}^{C}p(z_{i}|\mathcal{X},\mathcal{L},\mathcal{P}^{\prime})ln\bigg{(}p(\mathcal{X},\mathcal{L}|z_{i})p(z_{i}|\mathcal{P})\bigg{)}$
$\displaystyle=$
$\displaystyle\sum_{i=1}^{C}p(z_{i}|\mathcal{X},\mathcal{L},\mathcal{P}^{\prime})ln(p(z_{i}|\mathcal{X},\mathcal{L},P))+\sum_{i=1}^{C}p(z_{i}|\mathcal{X},\mathcal{L},\mathcal{P}^{\prime})ln(p(\mathcal{X},\mathcal{L}|\mathcal{P}))$
$\displaystyle=$
$\displaystyle\sum_{i=1}^{C}p(z_{i}|\mathcal{X},\mathcal{L},\mathcal{P}^{\prime})ln(p(z_{i}|\mathcal{X},\mathcal{L},\mathcal{P}))+ln(p(\mathcal{X},\mathcal{L}|\mathcal{P}))$
(8)
Given $\mathcal{P}^{\prime}$, we can infer the latent variable $z$. Because
the hard assignment of $z$, the first term in the right hand side of Eq.
(Proof) assigns $z_{i}$ into one class. Note that
$p(z|\mathcal{X},\mathcal{L},\mathcal{P})ln(p(z|\mathcal{X},\mathcal{L},\mathcal{P}))$
is a monotonically increasing function, which means that by maximizing the
posterior likelihood $p(z|\mathcal{X},\mathcal{L},\mathcal{P})$ for each
instance, we can maximize Eq. (Proof) for the hard assignment case in Eq. (3).
Thus, the updating strategy in our algorithm is a special case of EM
algorithm, and it can converge into a local maximum as EM algorithm. Note that
in our implement, we infer the latent variable in cluster level. In other
words, to maximize $p(z_{i}|\mathcal{X},\mathcal{L},\mathcal{P}^{\prime})$, we
can include another latent variable $\pi_{j},j\in[1,K]$. In other words, we
need to maximize
$\sum_{j=1}^{K}p(z_{i},\pi_{j}|\mathcal{X},\mathcal{L},\mathcal{P}^{\prime})$,
which we can recursively determine the latent variable $\pi_{i}$ using an
embedded EM algorithm. Hence, our algorithm use two steps of EM algorithm, and
it can converge to a local maximum. Refer [23] for more details about the
convergence of EM algorithm.
### 3.4 Probabilistic understanding for the model
The latent SVM model [11, 1] propose to label instance $x_{i}$ in positive
bag, by maximize $p(z(x_{i})=1|x_{i})$, which is the optimal Bayes decision
rule.
Figure 1: Example of graphical representation only for one class (event).
$h_{i}^{j}$ is the hidden variable, $x_{i}^{j}$ is the observable input,
$y_{i}^{j}$ is the projection of $x_{i}^{j}$ in the subspace, $j\in[1,n_{i}]$,
and $n_{i}$ is the number of total training data for class $i$. The $K$
cluster centers $\mu_{i}=\\{\mu_{i}^{1},\mu_{i}^{2},...,\mu_{i}^{K}\\}$ is
determined by both $\pi_{i}$ and $w_{i}$. The graphical model of our method is
similar to GMM model in vertical. By adding $z_{i}$ into LDA, the graphical
model can handle latent variables.
Similarly, Eq. (3a) takes the same strategy as latent SVM to maximize a
posterior probability. Moreover, instead of only maximizing the $p(z=1|x)$, we
also maximize the joint probability $p(z=1,x)$, using the Bayes rule,
$p(z=1,x)=p(x)p(z=1|x)$. In this paper, we use Gaussian mixture model to
approximate the prior $p(x)$ in this generative model. We argue that to
maximize a joint probability is reasonable, because it considers both
discriminative (posterior probability) and representative (prior) property in
the video dataset. We give the graphical representation of our model in Fig.
(1).
## 4 Experiments and results
In this section, we perform experiments on various data sets to evaluate the
proposed techniques and compare it to other baseline methods. For all the
experiments, we set $T$ = 20 and $\beta=40$; and initialize uniformly weighted
$w_{i}$ and projection matrix $\mathcal{P}$ with LDA.
### 4.1 Classification on toy data sets
The MUSK data sets111www.cs.columbia.edu/ andrews/mil/datasets.html are the
benchmark data sets used in virtually all previous approaches and have been
described in detail in the landmark paper [8]. Both data sets, MUSK1 and
MUSK2, consist of descriptions of molecules using multiple low-energy
conformations. Each conformation is represented by a 166-dimensional feature
vector derived from surface properties. MUSK1 contains on average
approximately 6 conformation per molecule, while MUSK2 has on average more
than 60 conformations in each bag. The Corel data set consists three different
categories ( elephant , fox , tiger ) , and each instance is represented with
230 dimension features, characterized by color, texture and shape descriptors.
The data sets have 100 positive and 100 negative example images. The latter
have been randomly drawn from a pool of photos of other animals. We first use
PCA reduce its dimension into 40 for our method. For parameter setting, we set
$K$=3, $T$ = 20 and $N=4$ (namely the 4-Nearest-Neighbor (4NN) algorithm is
applied for classification). The averaged results of ten 10-fold cross-
validation runs are summarized in Table (1). We set LDA222we use the bag label
as the instance label to test its performance and MI-SVM as our baseline. We
can observe that both LFDA and JLFDA outperform MI-SVM on MUSK1 and Fox data
sets, while has comparative performance as MI-SVM on the others.
Data set | inst/Dim | MI-SVM | LDA | LFDA | JLFDA
---|---|---|---|---|---
MUSK1 | 476/166 | 77.9 | 70.4 | 81.4 | 87.1
MUSK2 | 6598/166 | 84.3 | 51.8 | 76.4 | 81.3
Elephant | 1391/230 | 81.4 | 70.5 | 74.5 | 79.0
Fox | 1320/230 | 57.8 | 53.5 | 61.5 | 59.5
Tiger | 1220/230 | 84.0 | 71.5 | 74.0 | 80.5
Table 1: Accuracy results for various methods on MUSK and Corel data sets. Our
approach outperform LDA on both datasets, and we get better result than MI-SVM
on MUSK1 and Fox data set.
### 4.2 Semantic keyframe extraction
We conduct experiments on the challenging TRECVID MED11
dataset333http://www.nist.gov/itl/iad/mig/med11.cfm. It contains five events:
attempting a board trick feeding an animal, landing a fish, wedding ceremony
and working on a woodworking project. All of five events consist of a number
of human actions, processes, and activities interacting with other people
and/or objects under different place and time. At this moment, we take 105
videos from 5 events for testing and the remaining 710 videos for training.
For parameters, we set $K=10$ and $N=10$. We learned the representative
clusters for each class, and then use them to find semantic frames in videos
with the same labels. Then we evaluation the semantic frames for each video
through human-factors analysis—the semantic keyframe extraction problem
demands a human-in-the-loop for evaluation. We explain our human factors
experiment in full detail in experiment setup. Our ultimate findings
demonstrate that our proposed latent FDA with prior model is most capable of
extraction semantically meaningful keyframes among latent FDA and competitive
baselines.
Video representation. For all videos, we extract HOG3D descriptors [16] every
25 frames (about sampling a frame per second). To represent videos using local
features we apply a bag-of-words model, using all detected points and a
codebook with 1000 elements.
Benchmark methods. We make use of SVM as the benchmark method in the
experiment. We take the one-vs-all strategy to train a linear MI-SVM
classifier using $SVM^{light}$ [15], which is very fast in linear time, for
each kind of event. Then we choose 10 frames for each video which are far from
the margin and close to the margin on positive side. For the frames chosen
farthest away from the margin, we refer it SVM(1), while for frames closest to
the margin we refer it SVM(2). We also randomly select 10 frames from each
video, and we refer it RAND in our experiments.
Experiment setup Ten highly motivated graduate students (range from 22 to 30
years’ old) served as subjects in all of the following human-in-the-loop
experiments. Each novel subject to the annotation-task paradigm underwent a
training process. Two of the authors gave a detailed description about the
dataset and problem, including its background, definition and its purpose. In
order to indicate what representative and discriminative means for each event,
the two authors showed videos for each kind of event to the subjects, and make
sure all subjects understand what semantic keyframes are. The training
procedure was terminated after the subject’s performance had stabilized. We
take a pairwise ranking strategy for our evaluation. We extract 10 frames per
video for 5 different methods (SVM(1), SVM(2), LFDA, JLFDA and RAND)
respectively. For each video, we had about 1000 image pairs for comparison. We
had developed an interface using Matlab to display two image pair and three
options (Yes, No and Equal) to compare an image pair each time. The students
are taught how to use the software; a trial requires them to give a ranking:
If the left is better than the right, then choose ’Yes’; if the right is
better than the left, choose ’No’. If the two image pair are same, then choose
’Equal’. The subjects are again informed that better means a better semantic
keyframe. The ten subjects each installed the software to their computers, and
conducted the image pair comparison independently. In order to speed up the
annotation process, the interface can randomly sample 200 pairs from the total
1000 image pairs for each video, and we also ask subjects to random choose 10
videos from the test dataset.
Experimental Results We have scores for each image pair. By sampling 10 videos
from each event, we at last had annotations of 104 videos. It means our
sampling videos got from 10 subjects almost cover all test data (105 videos).
Table (2) shows the win-loss matrix between five methods by counting the
pairwise comparison results on all 5 events. It shows that JLFDA and LFDA
always beat the three baseline methods. Furthermore, JLFDA is better than LFDA
because it considers a prior distribution from training data, which will help
JLFDA to find more representative frames. See Fig. (3) for keyframes extracted
with JLFDA.
Figure 2: Comparison of 5 methods for five events. Higher value, better
performance.
Method | Win-Loss matrix
---|---
JLFDA | FLDA | RAND | SVM(1) | SVM(2)
JLFDA | - | 3413 | 2274 | 2257 | 3758
LFDA | 2957 | - | 2309 | 2230 | 3554
RAND | 2111 | 2175 | - | 1861 | 2274
SVM(1) | 2088 | 2270 | 2010 | - | 2314
SVM(2) | 3232 | 3316 | 2113 | 2125 | -
Table 2: Win-Loss matrix for five methods. It represents how many times
methods in each row win methods in column.
We compared the five methods on the basis of Condorcet voting method. We treat
’Yes’, ’No’ and ’Equal’ are voters for each method in the image pairwise
comparison process. If ’Yes’, we cast one ballot to the left method; else if
’No’, we add a ballot to the right method; else do nothing to the two methods.
Fig. (2) shows ballots for each method on each event. It demonstrates our
method JLFDA always beat other methods, except for E004 dataset. We also
compared the five methods based on Elo rating system. For each video, we
ranked the five methods according to Elo ranking system. Then, we counted the
number of No.1 methods in each event. The results in Table (3) show that our
method is better than others, except E004. Such results based on Elo ranking
is consistent with Condorcet ranking method in Fig. (2). E004 is the wedding
ceremony event and our method is consistently outperformed by the SVM baseline
method. We believe this is due to the distinct nature of the E004 videos in
which the video scene context itself distinguishes it from the other four
events (the wedding ceremonies typically have very many people and are
inside). Hence the discriminative component of the methods are taking over,
and the SVM is able to outperform the Fisher discriminant. This effect seems
more likely due to the nature of the five events in the data set than the
proposed method intrinsically.
Method | the number of No.1 method in each event
---|---
E001 | E002 | E003 | E004 | E005
JLFDA | 6 | 7 | 7 | 3 | 7
LFDA | 6 | 4 | 4 | 5 | 1
SVM(1) | 4 | 4 | 4 | 2 | 4
SVM(2) | 6 | 3 | 1 | 7 | 6
RAND | 2 | 2 | 4 | 3 | 2
Table 3: For each video, we ranked the five methods according to Elo ranking
system. Then, we counted the number of No.1 method one video level in each
event. For example, E002 has total 20 videos, and JLFDA has rank first on 7
videos, while RAND has rank first on only two videos. Higher value, better
results. It demonstrates that our method is more capable at extracting
semantically meaningful keyframes.
---
Figure 3: Sample keyframes from the first five events (each row from top to
down): (a) snowboard trick, (b) feeding animal, (c) fishing, (d) marriage
ceremony, and (e) wood making. Each row indicates sample results from the same
videos for each event. It shows that our method can extract the representative
and discriminative images for each kind of events. In other words, we can
decide what’s happen when we scan the images.
## 5 Conclusion
In this paper, we have presented a latent Fisher discriminant analysis model,
which combines the latent variable inference and dimension reduction in an
unified framework. Ongoing work will extend the kernel trick into the model.
We test our method on classification and semantic keyframe extraction problem,
and yield quite competitive results. To the best of our knowledge, this is the
first paper to study the extraction of semantically representative and
discriminative keyframes—most keyframe extraction and video summarization
focus on representation summaries rather than jointly representative and
discriminative ones. We have conducted a thorough ranking-based human factors
experiment for the semantic keyframe extraction on the challenging TRECVID
MED11 data set and found that our proposed methods are able to consistently
outperform competitive baselines.
## References
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|
arxiv-papers
| 2013-09-21T03:42:04 |
2024-09-04T02:49:51.249837
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Gang Chen",
"submitter": "Gang Chen",
"url": "https://arxiv.org/abs/1309.5427"
}
|
1309.5459
|
# Advances in computer architecture
Irfan Uddin Irfan Uddin
University of Amsterdam, The Netherlands
[email protected]
###### Abstract
In the past, efforts were taken to improve the performance of a processor via
frequency scaling. However, industry has reached the limits of increasing the
frequency and therefore concurrent execution of instructions on multiple cores
seems the only possible option. It is not enough to provide concurrent
execution by the hardware, software also have to introduce concurrency in
order to exploit the parallelism.
###### Contents
1. 1 Introduction
2. 2 Performance via frequency scaling
3. 3 Performance via automatic exploitation of ILP
4. 4 Performance via concurrency in hardware
5. 5 Performance via explicit concurrency in software
6. 6 Modern multi-core and many-core systems
7. 7 Conclusion
## 1 Introduction
In 1965 Intel’s co-founder Gorden Moore presented the now famous _Moore’s law_
stating that: _The number of transistors on a chip will roughly double every
18 months_. Later in 2000 the 18 months were changed to two years but the law
is still followed by electrical engineers. To quote Prof. Yale N. Patt:
_Moore’s law is not the law of physics, it is merely the psychology of
electrical engineers what they can do with silicon. The law is alive only
because it is used by industries to compete with each other_ [1].
Computer architects have consensus that in order to increase the performance
of the processor will likely require more transistors on the chip, to support
ever increasing software needs. While Moore’s law seems to enable the increase
in the number of transistors, it is not yet clear how these transistors can be
used to achieve the desired improvements. In this paper we summarize this
argument and promotes multi-core design as a promising approach. We present
some technological competition to the Microgrid.
The rest of the paper is organized as follows. In Section 2 we describe the
performance improvement in architecture by frequency scaling. We explain the
automatic exploitation of ILP to increase performance in Section 3 and the use
of concurrency in hardware to improve performance in Section 4. We show the
explicit concurrency in software to achieve performance in Section 5. We give
the details of modern many-core systems in Section 6 and conclude the paper in
Section 7.
## 2 Performance via frequency scaling
Since the advent of microprocessors, computer architects have always been in
the quest to achieve the best performance by increasing the frequency of a
single chip. To main focus is to achieve the lowest possible latency in the
execution of individual instructions and highest possible throughput in the
execution of the program. In order to achieve this improvement in performance
via frequency scaling, there are two possibilities; reduce the line width of
the technology and hence increase the switching speed of the transistors or
reduce the amount of work done in a cycle by pipelining instruction execution.
### 2.1 Power consumption
In order to achieve a higher frequency in a given technology a higher voltage
is required which means more power is consumed. However, power consumption is
at odds with ”Green computing” [23]. Also as the power consumption increases,
all kinds of problems come up: atom migration in gates, change of electronic
properties, wear and tear due to thermal expansion of the material itself,
etc. Furthermore, devices are getting portable, and high voltage means
constant access to power which is simply not practical in a portable
environment.
### 2.2 Dissipation of heat
The higher power consumed by a transistor is dissipated from the chip in the
form of heat which must be removed from the chip to avoid damaging it, as
silicon simply can not take heat. Increasing the frequency of the transistor
will increase the quantity of heat, and decreasing the size of the transistor
(decreases the energy required to switch a transistor) will generally increase
the density of heat. It is expensive to provide cooling to a heated processor
and requires extra space to install fans. Portable devices with this much heat
can burn the body of the user. There is also the fact that batteries store
limited energy and hence we want to minimise the power consumption in order to
maximise lifetime.
### 2.3 Delay in wires
The frequency of transistors can easily be improved, but wires are slow and
wire delay becomes slower for a given amount of power dissipation. A wire is
passive and can be made fast but then it will require more energy. It is
required to look at the relative speed when driving a wire by a minimum sized
transistor. Therefore the transaction from one part of the chip to another
part is limited by wires i.e. the frequency of transistors is not limited by
the number of transistors that can be manufactured on a chip, but actually
limited by the number of gates reachable in a cycle [25, 2]. Increasing the
frequency of the transistor means that there is a higher communication delay
than the computation delay.
In the past years the frequency of the processor has increased at the rate of
75% per year (not any more) and the frequency of the DRAM has increased at the
rate of just 7% per year [26]. Because of this difference, the DRAM latency
has become increasingly large compared to the processor’s cycle time which is
very high from processor’s point of view. This divergence in latency is also
known as ”Memory wall” [48]. To avoid the delay for DRAM access, concurrency
is the next logical step i.e. at the time the processor sends request to DRAM
and then has to wait for the request to complete, the processor should be able
to perform some other activities concurrent to the DRAM access.
## 3 Performance via automatic exploitation of ILP
An instruction stream has a combination of memory and computational
operations. When the memory operations are issued, there can be some
independent computational instructions that can be executed. But additional
hardware is required to find these independent instructions and schedule them
so that no data hazard occurs. This technique is called automatic exploitation
of Instruction Level Parallelism (ILP) [47]; and a number of approaches are
tried in this direction such as pipelining, out-of-order execution, branch
prediction and cache prefetching.
### 3.1 Pipelining
The function of each instruction (i.e. fetch, decode, execute etc.) can be
broken down in multiple sub-functions. Like a conveyor belt in a factory line,
these sub-functions can be pipelined in parallel where multiple instructions
coming one after another. Pipelining have enabled the simultaneous execution
of multiple instructions. The actual execution time of the individual
instruction remain the same, but the throughput of several instructions is
improved typically by $x$ times where $x$ is the length of the pipeline. e.g.
in Pentium machines the throughput of instructions was increased up to 20
times. In addition, pipelining also enabled the processor to have a higher
clock frequency because the operations executed at each pipeline stage is
simpler and therefore takes a shorter cycle time.
### 3.2 Out-of-order execution
The instruction stream has a limited number of independent instructions that
can be executed in parallel. Often, an instruction is dependent on a previous
operation (not only loads but also long latency operations like mul, floating
point operations etc.) and must therefore stall the pipeline until the load
operation completes. Out-of-order execution is used to allow instructions in
the pipeline to overtake those issued earlier in order to avoid stalling the
pipeline completely. However, hardware logic and energy is required for the
dependency analysis to determine if instructions can overtake each other.
Also, since instructions are executing out-of-order, a reorder buffer is
required for the in-order completion to provide the expected results [46].
Out-of-order execution introduces additional dependencies e.g. Write-After-
Read and Write-After-Write which are resolved by register renaming i.e.
allowing the hardware to use more registers than will fit in the instruction
format. However it increases the size of the register file, lengthens the
pipeline and therefore increases the branch penalties [18].
### 3.3 Branch prediction
The number of instructions between branches is typically very small (average
is less than 6 instructions [31]). Branch prediction is required to keep the
pipeline full; while the branch and instructions before it are executing, we
can fetch and decode instructions after the branch. If the branch was
predicted accurately we can already have the next instructions in the
pipeline, however if the branch was not predicted correctly it will result in
a pipeline bubble or many pipeline bubbles in more complex processors. In
addition, the effect of instructions must also be cancelled, which might
involve the roll back of side effects. It means that a huge number of cycles
and energy are lost on computing something that did not really participate in
the required computation of the program. The issue with branch prediction is,
if there are multiple tasks/processes the branch predictor has to predict
branches not in only one but multiple programs interleaved over time. To keep
the high accuracy across heterogeneous codes, the size of the branch predictor
must grow. Because of multitasking branch prediction is an expensive approach.
### 3.4 Cache prefetching
To avoid the delay of a cache miss, a processor can look at the memory access
pattern of the program and guess what cache line will be accessed next and
prefetch the line. If the guess is correct, the large delay required to access
the memory is avoided, however if it is not correct a large amount of energy
and memory bandwidth is wasted to fetch the data from the off-chip memory that
will not be used. Intel Itanium, Pentium 4 and Transmeta Crusoe are some
examples where cache prefetching is used.
### 3.5 Discussion
Traditional and super-scalar machines [20] implement most of the techniques
described above in order to achieve the best performance by the automatic
exploitation of ILP. However, because of the additional hardware the
microprocessors are getting complex, energy inefficient and not scalable [8].
Next to super-scalar machines, VLIW [11] machines have been introduced which
in general are more scalable than super-scalar machines but have a complicated
design and a complex code generation process. In addition the binary code for
a VLIW machine can not be used on the same VLIW machine with a different issue
width i.e. binary compatibility is not achieved [18].
Computer architects have tried hard to improve performance of the programs
through implicit ILP without changing anything in the software. Because of
implicit ILP, the programmers were not required to have a detailed knowledge
of the hardware. The programmer used to wait to buy until a new more powerful
machine is available in the market and then the software magically got faster.
However, the industry is reaching the limits of increasing the frequency of
processors to execute more instructions in one cycle [3].
It seems like the free lunch is over [33], programmers have to take
responsibility of parallelization in order to get more performance and
therefore they need to get familiar with the architecture and concurrent
execution model. Some automatic parallelization techniques are taking the
approach of abstracting away the architecture details [13], but can not
exploit the architectural resources fully as programmers can do manually with
the knowledge of the architecture. The introduction of concurrency in software
engineering, increases the complexity of the software, decreases the
productivity of the engineer, but increases the performance of the software.
## 4 Performance via concurrency in hardware
The increased number of transistors on a chip has enabled computer architects
to provide concurrent execution of independent instructions on multiple
execution units. However, the implicit concurrency that can be extracted
automatically in hardware is limited [3] i.e. parallelism is not possible
unless concurrency is introduced explicitly in programs.
## 5 Performance via explicit concurrency in software
Instead of using a single instruction stream and trying to improve the
performance of the stream by implicit parallelism, we can have multiple
instruction streams which can execute concurrently. These streams are called
”threads” and provide an independent flow of control within a process. A
thread is also called a ”light-weight” process, as it has its own stack, local
variables and program counter. Context switching in threads is much cheaper
than in processes. Industry has realized that writing parallel programs to
exploit concurrency of the hardware can not be avoided in future.
Introducing concurrency in programs is difficult, also debugging these
programs require a lot of effort from programmers. The order of execution is
known but the order of execution of threads is not known which makes the
behavior of multiple threads in a parallel program difficult to predict. The
synchronization between threads to ensure that shared state is accessed
atomically can make the concurrency difficult to define and can significantly
affect the performance of the program in multi-core systems. It is true that
threads in some form were quite commonly used in networking from 1980s e.g.
multiple connections to HTTP servers. In networking, threads represent
independent transaction and in applications, threads interact with each other
in producing results. Because of the inherent difficulty, parallel programming
never became a mainstream programming paradigm [34] in software engineering
community in the past.
Despite being difficult, parallel programing is the most desirable technology
in the current state-of-the-art multi-cores processors. The performance of an
application can be improved only from parallelization on contemporary
hardware. A number of programming libraries have been developed to provide
concurrency constructs to expose concurrency in programs. POSIX threads [12],
OpenMP [27], SystemC [14] etc. are some of the libraries that provide
constructs for creation, termination, synchronization and communication
between threads.
While parallelization is desirable, the management of threads in software by
the contemporary hardware is expensive. Typically 10-100 thousand cycles are
consumed in creation or synchronization of threads. There is also a cost of
context switching. Therefore fine-grained parallelism can not be achieved by
explicit concurrency in software. The next logical step is to introduce
concurrency at multiple levels; applications, operating systems and hardware.
We need threads in software, we need threads in hardware and we need the
management of threads in hardware. The concurrency at all levels will exploit
the maximum possible parallelism.
## 6 Modern multi-core and many-core systems
Since parallelization is the only practical solution in current technology and
the number of transistors on a single chip is growing [32], therefore we claim
that in the future there will be large number of cores on a single chip where
programmers have to put effort to write parallel programs with the knowledge
of the underlying architecture. In this section we describe some state-of-the-
art multi- and many- core systems. Some of these cores are available
commercially while others are mainly in the research domain and are comparable
to the Microgrid. There may exist some other many-core systems but these are
not discussed in this section.
### 6.1 Nvidia’s GPGPU
Nvidia’s GPGPU (General Purpose Graphical Processing Unit) [6] makes use of a
GPU as a co-processor to accelerate the execution in CPUs for general-purpose
computing. The acceleration happens by offloading some of the computationally
intensive and time consuming portions of the program from CPU to GPU. The rest
of the application still runs on the CPU. Therefore it is also known as
”heterogeneous” or ”hybrid” computing. A CPU generally consists of two, four
or eight cores, while the GPU consists of hundreds/thousands of smaller cores.
All these cores work together to crunch through the data in the application.
From the user’s point of view, the application runs faster because it is using
the massively parallel processing power of the GPU to boost performance. CUDA
(Computer Unified Device Architecture) is a parallel programming model for
GPU. It enables dramatic increase in computing performance by harnessing the
power of GPUs.
#### Discussion
GPU is actually a specialized accelerator that is connected to traditional
single- or multi-cores processor. Therefore a large amount of work is required
for the transaction of data between processor and GPUs. Programmers have to
divide the problem into coarse sub-problems that can be solved independently.
In addition the programmers have to explicitly manage the memory and
concurrency i.e. a complex model of the architecture is forced in the mind of
the programmer.
GPU architecture is based on the SIMD model and therefore can efficiently
execute SIMD based applications. The SIMD architectures are very inefficient
in branching, as each branch path must be executed sequentially. GPUs can
achieve a very high performance in embarrassingly parallel applications.
However applications with dense communication between threads (e.g. FFTs,
compression algorithms, sorting etc.) can not achieve a very high performance
compared to other multi-cores processors [21]. GPUs become really slow in
functional calls, e.g. FFT is an embarrassingly parallel application, but it
does not scale very well on GPUs because of function calls from the outer
loop.
### 6.2 Sun/Oracle’s UltraSPARC Tx
Sun Microsystems have introduced a RISC architecture named SPARC (Scalable
Processor ARChitecture) in 1987. Oracle then bought Sun Microsystems and they
together introduced UltraSPARC T1 microprocessor (code name Niagara) in 2005.
This continued as a series of UltraSPARC T2, UltraSPARC T3 and in 2011
UltraSPARC T4.
The UltraSPARC T4 [7] has a 16-stage integer pipeline, 11-stage floating point
pipeline and 2 issue width. It has a thread priority mechanism where one
thread can get preferential access to a core’s hardware to give increased
performance. It has 4 processors on a single die and each processor consists
of 8 cores therefore a total of 32 cores are available on a single chip. Each
core has 8 hardware threads i.e. 64 hardware threads on a processor and 256
hardware threads on the chip. Each core can switch between eight threads using
a modified LRU (Least Recently Used) algorithm for thread selection. Each core
is associated with 16KB of L1 I-cache and D-cache and 128KB of L2-cache. Eight
cores share 4MB L3-cache and the DDR is 1TB. Total transistors count is
approximately 855 millions. The frequency of every core can be changed in the
range of 2.85 and 3.0 GHz. The technology used is 40nm CMOS and the total die
size is $403mm^{2}$.
The UltraSPARC T4 processor has increased single-thread performance, while
maintaining the high multi-thread throughput performance, therefore single-
threaded applications can have an efficient execution. It automatically
switches to single-thread mode when only a single thread is active, dedicating
all resources to that thread’s execution. While software can activate up to
eight threads on each core at a time, hardware dynamically and seamlessly
allocates core resources such as instruction, data, L2-caches and TLBs, as
well as out-of-order execution resources such as the 128-entry re-order buffer
in the core. The cores provide sophisticated branch prediction and have the
features for prefetching instructions and data.
#### Discussion
The UltraSPARC is addressing internet servers or any large scale systems. It
generally addresses coarse-grained or embarrassingly parallel applications and
therefore disregard desktop computing and fine-grained parallelism. A library
is used to map software threads to hardware threads and it has an overhead of
creation and synchronization [43]. It is based on SMP model and has an
increased single thread performance. But it suffers from the scalability of
the bandwidth and power consumption of the interconnection between processor
and memory.
The UltraSPARC T4 cores are complex (16 stage pipeline) and have out-of-order
execution, branch prediction and cache prefetching, which are energy in-
efficient features. It has some inefficiency coming from the use of a huge
shared L2-cache which is necessary for server application (large number of
synchronizations around data) but has a cost in silicon i.e. heat, energy and
wiring complexity.
### 6.3 Tilera’s TILE64
Tilera has introduced TILE64 [29], based on MIMD model. It has 64 cores on a
chip using 90nm CMOS. These cores are fully functional, programmable and each
is capable of running its own operating system. A group of cores can also be
used to run as a symmetrical multi-processing operating system. Every core has
a frequency range of 600 to 900 MHz. The cores are in-order, three-way VLIW
issue width and implement a MIPS-derived VLIW instruction set. The pipeline
has many (more than 6) stages. Each core has 32 general-purpose registers and
three functional units: two integer arithmetic logic units and a load-store
unit. Every core has L1-cache and L2-cache as well as a distributed L3-cache.
Tilera’s architecture eliminates the on-chip bus interconnection by placing a
communications switch on each core and arranging them in a grid fashion on the
chip to create an efficient two-dimensional traffic system for packets. This
technology is named as intelligent mesh (iMesh). iMesh is similar to mesh
network used in Intel’s SCC or NoC in embedded systems, with the innovation
that the flow of the messages in the mesh network can dynamically be adapted
based on the load of the network.
Tilera’s Multicore Development Environment (MDE) provides programming
framework that allows developers to scale their applications to large-scale
multicore systems. It has enable the standard tools such as gcc, gdb and
oprofile, to be multi-core aware so that the developer can easily debug,
profile and optimize code running across dozens of cores.
#### Discussion
Tilera’s architecture does not include hardware support for floating point
operation and therefore is not suitable for scientific computing. It mainly
targets embedded applications e.g. video encoding/decoding and network packet
processing. It is programmed in such a way that requires using registers for
communication by the programmer. Which means there is more responsibility on
the part of the programmer or compiler in creating threads. In a way
programmers need to understand the architecture in detail in order to program
it. Also the mapping of threads in the program to the hardware requires a
software library and therefore the cost of creation, synchronization and
mapping of software threads to hardware can not be avoided.
### 6.4 Intel’s SCC
Intel’s SCC (Single-chip Cloud Computer) [15] is an experimental many-cores
research platform designed to address hardware and software challenges by
industry and academic institutions in the tera-scale project [28]. It consists
of 48 Pentium 1 cores connected in a mesh network and on-chip message passing
network for inter-thread communication. The cores are relatively simple but
fully functional general-purpose cores. There is no hardware cache coherency
protocol, which allowed Intel to place 48 cores on a chip using CMOS 45nm
technology. It does not come as a stand-alone computer and a management PC
(MCPC) is used to run the applications on the chip.
Intel’s SCC has fine-grained power management where software applications are
given control to turn cores on and off or to change their performance levels,
continuously adapting to use the minimum energy needed at a given moment. It
can run all 48 cores at one time over a range of 25W to 125W and selectively
vary the voltage and frequency of the mesh network as well as a group of
cores. Each tile (2 cores) can have its own frequency, and groupings of four
tiles (8 cores) can each run at their own voltage. Every core uses the
mainstream x86 (CISC) instruction set. The Linux operating system is available
for the chip, as well as gcc and Fortran compilers. A small library RCC is
used for the communication between cores.
#### Discussion
Intel’s SCC is a prototype which is designed for studying the parallel
programming paradigm in general-purpose computers. Therefore it is not really
a commercial product to be used for mainstream computing. It mainly addresses
coarse-grained parallelism, and may not achieve a high performance improvement
in fine-grained applications. The absence of hardware cache coherency protocol
places more responsibility on the programmer and hence requires more effort
from the programmers to manage the coherency of the caches. The Pentium 1 core
is actually single-threaded machine and therefore can not achieve latency
tolerance in long latency operations.
### 6.5 Microgrid
The Microgrid [19, 4, 16] is a general-purpose, many-core architecture
developed at the University of Amsterdam which implements hardware multi-
threading using data flow scheduling and a concurrency management protocol in
hardware to create and synchronize threads within and across the cores on
chip. The suggested concurrent programming model for this chip is based on
fork-join constructs, where each created thread can define further concurrency
hierarchically. This model is called the microthreading model and is also
applicable to current multi-core architectures using a library of the
concurrency constructs called _svp-ptl_ [45] built on top of pThreads. In our
work, we focus on a specific implementation of the microthreaded architecture
where each core contains a single issue, in-order RISC pipeline with an ISA
similar to DEC/Alpha, and all cores are connected to an on-chip distributed
memory network [17, 5]. Each core implements the concurrency constructs in its
instruction set and is able to support hundreds of threads and their contexts,
called microthreads and tens of families (i.e. ordered collections of
identical microthreads) simultaneously.
A number of tools and simulators are added to the designer’s toolbox and used
for the evaluation of the Microgrid from different perspective. The compiler
for the Microgrid [22] can generate binary for different implementations of
the Microgrid. We have software libraries that provide the run-time systems
for the microthreading model on the shared memory SMP machines and referred as
_svp-ptl_ [45] and distributed memory for clusters/grids and are referred as
Hydra [24] and _dsvp-ptl_ [44] The SL compiler can generate binary for UTLEON3
[9, 10], MGSim [5, 30] and HLSim [38, 39, 40, 37, 41, 42, 35, 36].
## 7 Conclusion
The psychology of electrical engineers is that they can double the number of
transistors on a single chip every second year, which has enabled computer
architects to design more complex microprocessors. A number of approaches were
tried to achieve improvements in the throughput of the program implicitly. But
industry has reached the limits of implicit improvement in performance, and
multi-core designs seem to be promising approaches to achieve performance
explicitly. However, concurrency in hardware alone can not improve the
performance of the program unless concurrency is also exposed in software.
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|
arxiv-papers
| 2013-09-21T10:33:36 |
2024-09-04T02:49:51.259126
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Irfan Uddin",
"submitter": "Irfan Uddin",
"url": "https://arxiv.org/abs/1309.5459"
}
|
1309.5461
|
# Linear kernels for $k$-tuple and liar’s domination
in bounded genus graphs
Arijit Bishnu
ACM Unit
Indian Statistical Institute
Kolkata, INDIA
[email protected] Arijit Ghosh 111 Supported by the Indo-German Max Planck
Center for Computer Science (IMPECS).
D1: Algorithms & Complexity
Max-Planck-Institut für Informatik
Saarbrücken, Germany
[email protected] Part of this work was done when the author was a visiting
scientist in ACM Unit, Indian Statistical Institute, Kolkata. Subhabrata Paul
ACM Unit
Indian Statistical Institute
Kolkata, INDIA
[email protected]
###### Abstract
A set $D\subseteq V$ is called a _$k$ -tuple dominating set_ of a graph
$G=(V,E)$ if $\left|N_{G}[v]\cap D\right|\geq k$ for all $v\in V$, where
$N_{G}[v]$ denotes the closed neighborhood of $v$. A set $D\subseteq V$ is
called a _liar’s dominating set_ of a graph $G=(V,E)$ if
* (i) $\left|N_{G}[v]\cap D\right|\geq 2$ for all $v\in V$, and
* (ii) for every pair of distinct vertices $u,v\in V$, $\left|(N_{G}[u]\cup N_{G}[v])\cap D\right|\geq 3$.
Given a graph $G$, the decision versions of $k$-Tuple Domination Problem and
the Liar’s Domination Problem are to check whether there exists a $k$-tuple
dominating set and a liar’s dominating set of $G$ of a given cardinality,
respectively. These two problems are known to be NP-complete [LC03, Sla09]. In
this paper, we study the parameterized complexity of these problems. We show
that the $k$-Tuple Domination Problem and the Liar’s Domination Problem are
W[2]-hard for general graphs but they admit linear kernels for graphs with
bounded genus.
##### Keywords.
$k$-tuple domination, liar’s domination, planar graphs, bounded genus graphs,
kernelization, and W[2]-hard
## 1 Introduction
Let $G=(V,E)$ be a graph. For a vertex $v\in V$, let $N_{G}(v)=\\{u\in V|uv\in
E\\}$ and $N_{G}[v]=N_{G}(v)\cup\\{v\\}$ denote the open and closed
neighborhoods of $v$, respectively. A set $D\subseteq V$ is called a
_dominating set_ of a graph $G=(V,E)$ if $|N_{G}[v]\cap D|\geq 1$ for all
$v\in V$. The _domination number_ of a graph $G$, denoted by $\gamma(G)$, is
the minimum cardinality of a dominating set of $G$. The concept of domination
has been well studied. Depending upon various applications, different
variations of domination have appeared in the literature [HHS98a, HHS98b].
Among different variations of domination, $k$-tuple domination and liar’s
domination are two important and well studied type of domination [HH00, LC02,
LC03, RS09, Sla09]. A set $D\subseteq V$ is called a _$k$ -tuple dominating
set_ of a graph $G=(V,E)$ if each vertex $v\in V$ is dominated by at least $k$
number of vertices in $D$, that is, $|N_{G}[v]\cap D|\geq k$ for all $v\in V$.
The concept of $k$-tuple domination in graphs was introduced in [HH00]. For
$k=2$ and $3$, it is called _double domination_ and _triple domination_
respectively. The _$k$ -tuple domination number_ of a graph $G$, denoted by
$\gamma_{k}(G)$, is the minimum cardinality of a $k$-tuple dominating set of
$G$. It is a simple observation that for the existence of a $k$-tuple
dominating set, we need $\delta(G)\geq k-1$, where $\delta(G)$ is the minimum
degree of $G$. On the other hand, liar’s domination is a new variation of
domination and was introduced in 2009 by Slater [Sla09]. A set $D\subseteq V$
is called a _liar’s dominating set_ of a graph $G=(V,E)$ if the following two
conditions are met:
condition (i)
$|N_{G}[v]\cap D|\geq 2$ for all $v\in V$
condition (ii)
for every pair of distinct vertices $u,v\in V$, $|(N_{G}[u]\cup N_{G}[v])\cap
D|\geq 3$
In a network guarding scenario, if sentinels are placed in the vertices of the
dominating set, then the graph (network) is guarded. Consider the situation
where a single sentinel is unreliable or lies and we do not know the exact
sentinel that lies. We then need a liar’s dominating set to guard the network.
The _liar’s domination number_ of a graph $G$, denoted by $\gamma_{LR}(G)$, is
the minimum cardinality of a liar’s dominating set of $G$. Formally, the
decision versions of $k$-Tuple Domination Problem and Liar’s Domination
Problem are defined as follows.
$k$-Tuple Domination Problem
_Instance:_ A graph $G=(V,E)$ and a nonnegative integer $p$.
_Question:_ Does there exist a $k$-tuple dominating set of cardinality at most
$p$?
Liar’s Domination Problem
_Instance:_ A graph $G=(V,E)$ and a nonnegative integer $p$.
_Question:_ Does there exist a liar’s dominating set of cardinality at most
$p$?
Note that, every liar’s dominating set is a double dominating set and every
triple dominating set is a liar’s dominating set. Hence, liar’s domination
number lies between double and triple domination number, that is,
$\gamma_{2}(G)\leq\gamma_{LR}(G)\leq\gamma_{3}(G)$.
The rest of the paper is organized as follows. Section $2$ introduces some
pertinent definitions and preliminary results that are used in the rest of the
paper and a brief review on the progress in the study of parametrization for
domination problems. Section $3$ deals with the hardness results of both
$k$-tuple domination problem and liar’s domination problem. In Section $4$, we
show that both $k$-Tuple Domination Problem and Liar’s Domination Problem
admit linear kernel in planar graphs. In Section $5$, we extend the results
for bounded genus graphs. Finally, Section $6$ concludes the paper.
## 2 Preliminaries
Let $G=(V,E)$ be a graph. Let $G[S]$, $S\subseteq V$ denote the induced
subgraph of $G$ on the vertex set $S$. The _distance_ between two vertices $u$
and $v$ in a graph $G$ is the number of edges in a shortest path connecting
them and is denoted as $d_{G}(u,v)$. The degree of a vertex $v\in V(G)$,
denoted by $deg_{G}(v)$, is the number of neighbors of $v$.
### 2.1 Graphs on surfaces
In this subsection, we recall some basic facts about graphs on surfaces
following the discussion in [FT04]. The readers are referred to [MT01] for
more details. A _surface_ $\Sigma$ is a compact $2$-manifold without boundary.
Let $\Sigma_{0}$ denote the sphere $\\{(x,y,z)|~{}x^{2}+y^{2}+z^{2}=1\\}$. A
_line_ and _O-arc_ are subsets of $\Sigma$ that are homeomorphic to $[0,1]$
and a circle respectively. A subset of $\Sigma$ meeting the drawing only in
vertices of $G$ is called _$G$ -normal_. If an O-arc is G-normal, then it is
called a _noose_. The length of a noose is the number of its vertices. The
_representativity_ of $G$ embedded in $\Sigma\neq\Sigma_{0}$ is the smallest
length of a non-contractible noose in $\Sigma$ and it is denoted by _rep
$(G)$_.
The classification theorem for surfaces states that, any surface $\Sigma$ is
homeomorphic to either a surface $\Sigma^{h}$ which is obtained from a sphere
by adding $h$ handles (orientable surface), or a surface $\Sigma^{k}$ which is
obtained from a sphere by adding $k$ crosscaps (non-orientable surface)
[MT01]. The _Euler genus_ of a non-orientable surface $\Sigma$, denoted by _eg
$(\Sigma)$_, is the number of crosscaps $k$ such that $\Sigma\cong\Sigma^{k}$
and for an orientable surface, _eg $(\Sigma)$_ is twice the number of handles
$h$ such that $\Sigma\cong\Sigma^{h}$. Given a graph $G$, Euler genus of $G$,
denoted by eg$(G)$, is the minimum eg$(\Sigma)$, where $\Sigma$ is a surface
in which $G$ can be embedded. The _Euler characteristic_ of a surface $\Sigma$
is defined as $\chi(\Sigma)=2-\textup{eg}(\Sigma)$. For a graph $G$, $\chi(G)$
denotes the largest number $t$ for which $G$ can be embedded on a surface
$\Sigma$ with $\chi(\Sigma)=t$. Let $G=(V,E)$ be a $2$-cell embedded graph in
$\Sigma$, that is, all the faces of $G$ is homeomorphic to an open disk. If
$F$ is the set of all faces, then _Euler’s formula_ tells that
$V-E+F=\chi(\Sigma)=2-\textup{eg}(\Sigma)$.
Next we define a process called _cutting along a noose $N$_. Although the
formal defi is given in [MT01], we follow a more intuitive defi given in
[FT04]. Let $N$ be a noose in a $\Sigma$-embedded graph $G=(V,E)$. Suppose for
any $v\in N\cap V$, there exists an open disk $\Delta$ such that $\Delta$
contains $v$ and for every edge $e$ adjacent to $v$, $e\cap\Delta$ is
connected. We also assume that $\Delta-N$ has two connected components
$\Delta_{1}$ and $\Delta_{2}$. Thus we can define partition of
$N_{G}(v)=N_{G}^{1}(v)\cup N_{G}^{2}(v)$, where $N_{G}^{1}(v)=\\{u\in
N_{G}(v)|uv\cap\Delta_{1}\neq\emptyset\\}$ and $N_{G}^{2}(v)=\\{u\in
N_{G}(v)|uv\cap\Delta_{2}\neq\emptyset\\}$. Now for each $v\in N\cap V$ we do
the following:
1. 1.
remove $v$ and its incident edges
2. 2.
introduce two new vertices $v^{1},v^{2}$ and
3. 3.
connect $v^{i}$ with the vertices in $N_{G}^{i}$, $i=1,2$.
The resulting graph $\mathcal{G}$ is obtained from $\Sigma$-embedded graph $G$
by cutting along $N$. The following lemma is very useful in the proofs by
induction on the genus.
###### Lemma 1
[FT04] Let $G$ be a $\Sigma$-embedded graph and let $\mathcal{G}$ be a graph
obtained from $G$ by cutting along a non-contractible noose $N$. Then one of
the following holds
* •
$\mathcal{G}$ is the disjoint union of graphs $G_{1}$ and $G_{2}$ that can be
embedded in surfaces $\Sigma_{1}$ and $\Sigma_{2}$ such that
$eg(\Sigma)=eg(\Sigma_{1})+eg(\Sigma_{2})$ and $eg(\Sigma_{i})>0,i=1,2$.
* •
$\mathcal{G}$ can be embedded in a surface with Euler genus strictly smaller
than eg$(\Sigma)$.
A _planar graph_ $G=(V,E)$ is a graph that can be embedded in the plane. We
term such an embedding as a _plane graph_.
### 2.2 Parameterization and domination
A _parameterized problem_ is a language
$L\subseteq\Sigma^{*}\times\mathds{N}$, where $\Sigma^{*}$ denotes the set of
all finite strings over a finite alphabet $\Sigma$. A parameterized problem
$L$ is _fixed-parameter tractable_ if the question “$(x,p)\in L$” can be
decided in time $f(p)\cdot|x|^{O(1)}$, where $f$ is a computable function on
nonnegative integers, $x$ is the instance of the problem and $p$ is the
parameter. The corresponding complexity class is called FPT. Next we define a
reducibility concept between two parameterized problems.
###### Definition 2
[DF99, Nie06] Let $L,L^{\prime}\subseteq\Sigma^{*}\times\mathds{N}$ be two
parameterized problems. We say that $L$ reduces to $L^{\prime}$ by a _standard
parameterized m-reduction_ if there are functions $p\mapsto p^{\prime}$ and
$p\mapsto p^{\prime\prime}$ from $\mathds{N}$ to $\mathds{N}$ and a function
$(x,p)\mapsto x^{\prime}$ from $\Sigma^{*}\times\mathds{N}$ to $\Sigma^{*}$
such that
1. 1.
$(x,p)\mapsto x^{\prime}$ is computable in time $p^{\prime\prime}|x|^{c}$ for
some constant $c$ and
2. 2.
$(x,p)\in L$ if and only if $(x^{\prime},p^{\prime})\in L^{\prime}$.
A parameterized problem is in the class W[i], if every instance $(x,p)$ can be
transformed (in fpt-time) to a combinatorial circuit that has height at most
$i$, such that $(x,p)\in L$ if and only if there is a satisfying assignment to
the inputs, which assigns $1$ to at most $p$ inputs. A problem $L$ is said to
be W[i]_-hard_ if there exists a standard parameterized m-reduction from all
the problems in W[i] to $L$ and in addition, if the problem is in W[i], then
it is called W[i]_-complete_.
Next we define the reduction to problem kernel, also simply referred to as
_kernelization_.
###### Definition 3
[Nie06] Let $L$ be a parameterized problem. By reduction to problem kernel, we
mean to replace instance $I$ and the parameter $p$ of $L$ by a “reduced”
instance $I^{\prime}$ and by another parameter $p^{\prime}$ in polynomial time
such that
* •
$p^{\prime}\leq c\cdot p$, where $c$ is a constant,
* •
$I^{\prime}\leq g(p)$, where $g$ is a function that depends only on $p$, and
* •
$(I,p)\in L$ if and only if $(I^{\prime},p^{\prime})\in L$.
The reduced instance $I^{\prime}$ is called the _problem kernel_ and the size
of the problem kernel is said to be bounded by $g(p)$.
In parameterized complexity, domination and its variations are well studied
problems. The decision version of domination problem is W[2]-complete for
general graphs [DF99]. But this problem is FPT when restricted to planar
graphs [AFN04a] though it is still NP-complete for this graph class [GJ79].
Furthermore, for bounded genus graphs, which is a super class of planar
graphs, domination problem remains FPT [FT04]. It was proved that dominating
set problem possesses a linear kernel in planar graphs [AFN04a] and in bounded
genus graphs [FT04]. Also domination problem admits polynomial kernel on
graphs excluding a fixed graph $H$ as a minor [Gut09] and on $d$-degenerated
graphs [PRS12]. A search tree based algorithm for domination problem on planar
graphs, which runs in $O(8^{p}n)$ time, is proposed in [AFF+05]. For bounded
genus graphs, similar search tree based algorithm is proposed in [EFF04] and
has a time complexity of $O((4g+40)^{p}n^{2})$, where $g$ is the genus of the
graph. Algorithms with running time of $O(c^{\sqrt{p}}n)$ for domination
problem on planar graphs have been devised in [ABF+02, AFN04b, FT03, FT04].
Like domination problem, $k$-Tuple Domination Problem and Liar’s Domination
Problem are both NP-complete [LC03, Sla09] for general graphs. However, these
problems have been polynomially solved for different graph classes [LC02,
LC03, PP13a, PP13b]. But for planar graphs and hence for graphs with bounded
genus, $k$-Tuple Domination Problem remains NP-complete [LC08]. In [Sla09],
though the NP-completeness proof is given for general graphs, it can be
verified that using the same construction one can find the NP-completeness of
Liar’s Domination Problem in planar graphs, see Lemma 26 in Appendix 7.
Some generalization of classical domination problem have been studied in the
literature from parameterized point of view. Among those problems,
$k$-dominating threshold set problem, $[\sigma,\rho]$-dominating set problem
(also known as generalized domination) are generalized version of $k$-tuple
dominating set problem. In [GV08], it is proved that $k$-dominating threshold
set problem is FPT in $d$-degenerated graphs. $[\sigma,\rho]$-domination is
studied in [CP14, TP97, vRBR09]. A set $D$ of vertices of a graph $G$ is
$[\sigma,\rho]$-dominating set if for any $v\in D,|N(v)\cap D|\in\sigma$ and
for any $v\notin D,|N(v)\cap D|\in\rho$ for any two sets $\sigma$ and $\rho$.
It is known that $[\sigma,\rho]$-domination is FPT when parameterized by
treewidth [vRBR09]. By Theorem 32 of [ABF+02], it follows that $k$-tuple
domination is FPT on planar graphs. But there is no explicit kernel for both
$k$-tuple domination and liar’s domination problem in the literature.
There have been successful efforts in developing meta-theorems like the
celebrated Courcelle’s theorem [Cou92] which states that all graph properties
definable in monadic second order logic can be decided in linear time on
graphs of bounded tree-width. This also implies FPT algorithms for bounded
tree-width graph for these problems. In case of kernelization in bounded genus
graphs, Bodlaender et al. give two meta-theorems [BFL+09]. The first theorem
says that all problems expressible in counting monadic second order (CMSO)
logic and satisfying a coverability property admit a polynomial kernel on
graphs of bounded genus and the second theorem says that all problems that
have a finite integer index and satisfy a weaker coverability property admit a
linear kernel on graphs of bounded genus. It is easy to see that both
$k$-tuple and liar’s domination problems can be expressed in CMSO logic. Let
$G=(V,E)$ be an instance of a graph problem $\Pi$ such that $G$ is embeddable
in a surface of Euler genus at most $r$. The basic idea of quasi-coverable
property for $\Pi$ is that there exists a set $S\subseteq V$ satisfying the
conditions of $\Pi$ such that the tree-width of $G\setminus R^{r}_{G}(S)$ is
at most $r$ where $R^{r}_{G}(S)$ is a special type of reachability set from
$S$. In domination type of problems, this reachability set is actually the
whole graph and hence these problems satisfy the quasi-coverable property. The
basic idea of strong monotonicity for a graph problem $\Pi$ is roughly as
follows: Let $\mathcal{F}_{i}$ be a class of graphs $G$ having a specific set
of vertices $S$ termed as the boundary of $G$ such that $|S|=i$. The glued
graph $G=G_{1}\oplus G_{2}$ of $G_{1}$ and $G_{2}$ is the graph which is
obtained by taking the disjoint union of $G_{1}$ and $G_{2}$ and joining $i$
edges between the vertices of the boundary sets. A problem $\Pi$ is said to
satisfy the strong monotonicity if for every boundaried graph
$G=(V,E)\in\mathcal{F}_{i}$, there exists a set $W\subseteq V$ of a specific
cardinality which satisfy the property of $\Pi$ such that for every boundaried
graph $G^{\prime}=(V^{\prime},E^{\prime})\in\mathcal{F}_{i}$ with a set
$W^{\prime}\subseteq V^{\prime}$, satisfying the property of $\Pi$, the vertex
set $W\cup W^{\prime}$ satisfies the property of $\Pi$ for the glued graph
$G=G\oplus G^{\prime}$. It can be verified easily that both $k$-tuple
domination and liar’s domination problems satisfy the strongly monotone
property. The strongly monotone property implies the finite integer index for
these problems. Hence, by the second meta-theorem in [BFL+09], both $k$-tuple
and liar’s domination problems admit linear kernels for graphs on bounded
genus. Though these meta-theorems provide simple criteria to decide whether a
problem admits a linear or polynomial kernel, finding a linear kernel with
reasonably small constants for a specific problem is a worthy topic of further
research [BFL+09]. In this paper, we have obtained linear kernels with small
constants for both the problems on bounded genus graphs. We have also proved
the W[2]-hardness for $k$-tuple and liar’s domination for general graphs.
## 3 Hardness results in general graphs
In this section, we show that $k$-tuple Domination Problem and Liar’s
Domination Problem are W[2]-hard. In [CP14], it is proved that
$[\sigma,\rho]$-domination problem for any recursive sets $\sigma$ and $\rho$
is W[2]-hard. This implies the hardness for $k$-tuple domination in general
graphs. But in this paper, we have come up with a simple W[2]-hardness proof
for $k$-tuple domination in general graphs. To prove this, we show standard
parameterized m-reductions from Domination Problem, which is known to be
W[2]-complete [DF99], to $k$-Tuple Domination Problem and Liar’s Domination
Problem, respectively.
###### Theorem 4
$k$-Tuple Domination Problem is $\mathsf{W[2]}$-hard.
###### Proof.
We show a standard parameterized m-reduction from Domination Problem to
$k$-Tuple Domination Problem. Let $<G=(V,E),p>$ be an instance of Domination
Problem. We construct an instance
$<G^{\prime}=(V^{\prime},E^{\prime}),p^{\prime}>$ of the $k$-Tuple Domination
Problem as follows: $V^{\prime}=V\cup V_{k}$ where
$V_{k}=\\{u_{1},u_{2},\ldots,u_{k}\\}$ and
$E^{\prime}=E\cup\\{v_{i}u_{j}|v_{i}\in V\mbox{ and }u_{j}\in V_{k}\setminus
u_{k}\\}\cup\\{u_{i}u_{j}|u_{i},u_{j}\in V_{k},i\not=j\\}$. Also set
$p^{\prime}=p+k$. The construction of $G^{\prime}$ from $G$ in case of triple
domination is illustrated in Figure 1.
Figure 1: Construction of $G^{\prime}$ from $G$ for triple domination
###### Claim 5
$G$ has a dominating set of size at most $p$ if and only if $G^{\prime}$ has a
$k$-tuple dominating set of size at most $p^{\prime}$.
###### Proof.
Let $D$ be a dominating set of $G$ of cardinality at most $p$ and
$D^{\prime}=D\cup V_{k}$. Each $v_{i}\in V$ is dominated by at least one
vertex from $D$ and by $k-1$ vertices from $V_{k}$. Each $u_{i}\in V_{k}$ is
dominated by $k$ vertices of $V_{k}$. Thus, $D^{\prime}$ is a $k$-tuple
dominating set of $G^{\prime}$ of cardinality at most $p^{\prime}$.
Conversely, let $D^{\prime}$ be a $k$-tuple dominating set of $G^{\prime}$ of
cardinality at most $p^{\prime}$. Note that each $k$-tuple dominating set
contains the set $V_{k}$ because to dominate $u_{k}$ by $k$ vertices we must
select all the vertices of $V_{k}$. Let $D=D^{\prime}\setminus V_{k}$. Clearly
$D\subseteq V$ and $|D|\leq p$. Now for each $v\in V$, $|N_{G}[v]\cap D|\geq
1$ because otherwise, there exists a vertex $v\in V$ such that
$|N_{G^{\prime}}[v]\cap D^{\prime}|=k-1$. This is a contradiction because
$D^{\prime}$ is a $k$-tuple dominating set of $G^{\prime}$. Thus $D$ is a
dominating set of $G$ of cardinality at most $p$.
Hence, $G$ has a dominating set of size at most $p$ if and only if
$G^{\prime}$ has a $k$-tuple dominating set of size at most $p^{\prime}$. ∎
Thus, $k$-Tuple Domination Problem is W[2]-hard. ∎
Next we show the W[2]-hardness of Liar’s Domination Problem.
###### Theorem 6
Liar’s Domination Problem is $\mathsf{W[2]}$-hard.
###### Proof.
We show a standard parameterized m-reduction from Domination Problem to Liar’s
Domination Problem. Let $<G=(V,E),p>$ be an instance of Domination Problem. We
construct an instance $<G^{\prime}=(V^{\prime},E^{\prime}),p^{\prime}>$ of the
Liar’s Domination Problem as follows:
$V^{\prime}=V\cup\\{u,u^{\prime},v,v^{\prime},w\\}$ and
$E^{\prime}=E\cup\\{v_{i}u|v_{i}\in V\\}\cup\\{v_{i}v|v_{i}\in
V\\}\cup\\{uu^{\prime},vv^{\prime},wu,wv\\}$. Also $p^{\prime}=p+4$. The
construction of $G^{\prime}$ from $G$ is illustrated in Figure 2.
Figure 2: Construction of $G^{\prime}$ from $G$
###### Claim 7
$G$ has a dominating set of size at most $p$ if and only if $G^{\prime}$ has a
liar’s dominating set of size at most $p^{\prime}$.
###### Proof.
Let $D$ be a dominating set of $G$ of cardinality at most $p$ and
$D^{\prime}=D\cup\\{u,u^{\prime},v,v^{\prime}\\}$. It is easy to verify that
for each vertex $x\in V^{\prime}$, $|N_{G^{\prime}}[x]\cap D^{\prime}|\geq 2$
and for every pair of vertices $x,y\in V^{\prime}$, $|(N_{G^{\prime}}[x]\cup
N_{G^{\prime}}[y])\cap D^{\prime}|\geq 3$. Hence $D^{\prime}$ is a liar’s
dominating set of $G^{\prime}$ of cardinality at most $p+4=p^{\prime}$.
Conversely, let $D^{\prime}$ be a liar’s dominating set of $G^{\prime}$ of
cardinality at most $p^{\prime}$. Each liar’s dominating set must contain the
set $\\{u,u^{\prime},v,v^{\prime}\\}$ because to doubly dominate $u^{\prime}$
and $v^{\prime}$ we must select the vertices
$\\{u,u^{\prime},v,v^{\prime}\\}$.
Let $X\subseteq V$ denote the set of vertices that are dominated by exactly
two vertices ($u$ and $v$) from $D^{\prime}$. We claim $|X|\leq 1$. If there
exists two such vertices $x,y\in V$, then $|(N_{G^{\prime}}[x]\cup
N_{G^{\prime}}[y])\cap D^{\prime}|=2$ which violates condition (ii) of liar’s
domination. We now deal with two cases:
$|X|=1$:
Let $X=\\{x\\}$. Here $|N_{G^{\prime}}[x]\cap D^{\prime}|=2$. This implies
$w\in D^{\prime}$, otherwise the pair $x$ and $w$ violates condition (ii) of
liar’s domination. We set
$D^{\prime\prime}=(D^{\prime}\setminus\\{w\\})\cup\\{x\\}$. $D^{\prime\prime}$
is also a liar’s dominating set of $G^{\prime}$ of cardinality at most
$p^{\prime}$. Note that all vertices in $V$ is triply dominated by
$D^{\prime\prime}$ and it does not contain $w$.
$|X|=0$:
In this case each vertex of $V$ is triply dominated by $D^{\prime}$. Now if
$w\notin D^{\prime}$, we are done. Otherwise the set
$D^{\prime}\setminus\\{w\\}$ forms a liar’s dominating set of $G^{\prime}$ of
cardinality at most $p^{\prime}$ such that each vertex of $V$ is triply
dominated by $D^{\prime}\setminus\\{w\\}$.
Hence without loss of generality, we assume that there is a liar’s dominating
set $D^{\prime}$ of $G^{\prime}$ of cardinality at most $p^{\prime}$ such that
every vertex in $V$ is triply dominated by $D^{\prime}$ and $w\notin
D^{\prime}$. Let $D=D^{\prime}\setminus\\{u,u^{\prime},v,v^{\prime}\\}$.
Clearly $D\subseteq V$ and $|D|\leq p$. Now for each $x\in V$, $|N_{G}[x]\cap
D|\geq 1$ because otherwise, there exists a vertex $x\in V$ such that for the
pair $x$ and $w$, condition (ii) of liar’s domination is violated. This is a
contradiction because $D^{\prime}$ is a liar’s dominating set of $G^{\prime}$.
Thus $D$ is a dominating set of $G$ of cardinality at most $p$.
Hence, $G$ has a dominating set of size at most $p$ if and only if
$G^{\prime}$ has a liar’s dominating set of size at most $p^{\prime}$. ∎
Thus, Liar’s Domination Problem is W[2]-hard. ∎
## 4 Linear kernels for planar graphs
Having seen that $k$-tuple and liar’s domination are W[2]-hard in general
graphs, we focus on planar graphs in this section and show that they are FPT.
### 4.1 Double domination
In this subsection we show that Double Domination Problem in planar graphs
possesses a linear kernel. Our proof technique uses the region decomposition
idea of Alber et al. [AFN04a]. First we describe the reduction rules for
kernelization.
#### 4.1.1 Reduction rule
Let $G=(V,E)$ be the instance for Double Domination Problem. Consider a pair
of vertices $u,v\in V$. Let $N_{G}(u,v)=N_{G}(u)\cap N_{G}(v)$. We partition
the vertices of $N_{G}(u,v)$ in to three parts as follows:
$\displaystyle N^{1}_{G}(u,v)$ $\displaystyle=$ $\displaystyle\\{x\in
N_{G}(u,v)|N_{G}(x)\setminus\\{N_{G}(u,v)\cup\\{u,v\\}\\}\neq\emptyset\\};$
$\displaystyle N^{2}_{G}(u,v)$ $\displaystyle=$ $\displaystyle\\{x\in
N_{G}(u,v)\setminus N^{1}_{G}(u,v)|N_{G}(x)\cap
N^{1}_{G}(u,v)\neq\emptyset\\};$ $\displaystyle N^{3}_{G}(u,v)$
$\displaystyle=$ $\displaystyle N_{G}(u,v)\setminus(N^{1}_{G}(u,v)\cup
N^{2}_{G}(u,v)).$
Reduction Rule: For every pair of distinct vertices $u,v\in V$, if
$N^{3}_{G}(u,v)\neq\emptyset$, then
* •
delete all the vertices of $N^{2}_{G}(u,v)$ and
* •
delete all vertices of $N^{3}_{G}(u,v)$ except one vertex.
###### Lemma 8
Let $G=(V,E)$ be a graph and $G^{\prime}=(V^{\prime},E^{\prime})$ be the
resulting graph after having applied the reduction rule to $G$. Then
$\gamma_{2}(G)=\gamma_{2}(G^{\prime})$.
###### Proof.
Let $u,v\in V$ such that $N^{3}_{G}(u,v)\neq\emptyset$. Now if
$N^{2}_{G}(u,v)=\emptyset$ and $|N^{3}_{G}(u,v)|=1$, then $G^{\prime}$ is same
as $G$. So, without loss of generality, assume that
$\left|N^{3}_{G}(u,v)\right|>1$ and $N^{2}_{G}(u,v)\neq\emptyset$. Note that a
vertex $x$ of $N^{3}_{G}(u,v)$ can be doubly dominated by any two vertices
from $N_{G}[x]\subseteq\\{N^{2}_{G}(u,v)\cup N^{3}_{G}(u,v)\cup\\{u,v\\}\\}$.
Again for any two vertices $x,y\in N^{2}_{G}(u,v)\cup
N^{3}_{G}(u,v)\cup\\{u,v\\}$, $N_{G}[x]\cup N_{G}[y]\subseteq N_{G}[u]\cup
N_{G}[v]$. This shows that we can double dominate each vertex of
$N^{3}_{G}(u,v)$ in an optimal way by selecting $u$ and $v$ only. This
selection of $u$ and $v$ was forced by the only vertex $w\in N^{3}_{G}(u,v)$
that remained in $G^{\prime}$. We claim that $G$ contains a minimum double
dominating set $D$ which does not contain any vertex from $N^{2}_{G}(u,v)\cup
N^{3}_{G}(u,v)$. First observe that there can not be three or more vertices
from $N^{2}_{G}(u,v)\cup N^{3}_{G}(u,v)$ in $D$. If it were, then we could
replace those three or more vertices by $u$ and $v$, thus contradicting the
minimality of $D$. Now for those two (or one) vertices from
$N^{2}_{G}(u,v)\cup N^{3}_{G}(u,v)$ in $D$, we can replace them by $u$ and
(or) $v$. Therefore, $G$ contains a minimum double dominating set $D$ which
does not contain any vertex from $N^{2}_{G}(u,v)\cup N^{3}_{G}(u,v)$. Clearly,
this set $D$ also forms a minimum double dominating set of $G^{\prime}$.
Hence, $\gamma_{2}(G)=\gamma_{2}(G^{\prime})$. ∎
In this reduction, for a pair of distinct vertices $u,v\in V$, we have
actually deleted at most $\min\\{deg_{G}(u),deg_{G}(v)\\}$ vertices. So, the
time taken is $\sum_{u,v\in V}\min\\{deg_{G}(u),deg_{G}(v)\\}$ for the whole
reduction process. Since for a planar graph $\sum_{v\in V}deg_{G}(v)=O(n)$,
where $n$ is the number of vertices, we have the following lemma.
###### Lemma 9
For a planar graph having $n$ vertices, the reduction rule can be carried out
in $O(n^{3})$ time.
#### 4.1.2 A linear kernel
In this subsection, we show that the reduction rule given in the previous
section yields a linear kernel for Double Domination Problem in planar graphs.
For this proof, first we find a “maximal region decomposition” of the vertices
$V^{\prime}$ of the reduced graph $G^{\prime}=(V^{\prime},E^{\prime})$ and
then we show that $|V^{\prime}|=O(\gamma_{2}(G^{\prime}))$. We start with some
definitions regarding maximal region decomposition following Alber et al.
[AFN04a].
###### Definition 10
Let $G=(V,E)$ be a plane graph. A closed subset of the plane is called a
region $R(u,v)$ between two vertices $u,v$ if the following properties are
met:
1. 1.
the boundary of $R(u,v)$ is formed by two simple paths $P$ and $Q$ between $u$
and $v$ of length at most two edges, and
2. 2.
all the vertices which are strictly inside the region $R(u,v)$ are from
$N_{G}(u)\cap N_{G}(v)$.
The definition of a region is slightly different from the definition given in
[AFN04a], where all the vertices which are strictly inside the region $R(u,v)$
are from $N_{G}(u)\cup N_{G}(v)$. Note that by the above definition, paths of
length one or two between $u$ and $v$ can form a region $R(u,v)$. For a region
$R=R(u,v)$, let $\partial(R)$ denote the boundary of $R$ and $V(R)$ denote the
vertices inside or on the boundary of $R$, i.e., $V(R)=\\{u\in V|~{}u~{}\mbox{
is inside }R\mbox{ or on }\partial(R)\\}$.
###### Definition 11
Let $G=(V,E)$ be a plane graph and $D\subseteq V$. A $D$-region decomposition
of $G$ is a set $\mathcal{R}$ of regions between pairs of vertices in $D$ such
that
1. 1.
for $R(u,v)\in\mathcal{R}$ no vertices of $D$ (except $u$ and $v$) lies in
$V(R(u,v))$, and
2. 2.
for two regions $R_{1},R_{2}\in\mathcal{R}$, $(R_{1}\cap
R_{2})\subseteq(\partial(R_{1})\cup\partial(R_{2}))$, i.e., they can intersect
only at the vertices on the boundary.
For a $D$-region decomposition $\mathcal{R}$, we define
$V(\mathcal{R})=\cup_{R\in\mathcal{R}}V(R)$. A $D$-region decomposition
$\mathcal{R}$ is called maximal if there is no region $R$ such that
$\mathcal{R^{\prime}}=\mathcal{R}\cup R$ is a $D$-region decomposition, where
$V(\mathcal{R})$ is a strict subset of $V(\mathcal{R^{\prime}})$.
First we observe an important property of a maximal $D$-region decomposition.
###### Lemma 12
Let $G=(V,E)$ be a plane graph with a double dominating set $D$ and let
$\mathcal{R}$ be a maximal $D$-region decomposition. Then $V=V(\mathcal{R})$.
###### Proof.
Let $v\in V$ be a vertex such that $v\notin V(\mathcal{R})$. There can be two
cases – $v\in D$ and $v\notin D$. First, let us assume that $v\in D$. Since
$D$ is a double dominating set of $G$, there exists another vertex $x\in D$
such that $vx\in E$. Now, the path $P=(x,v)$ forms a region $R$. Clearly
$\mathcal{R}\cup R$ forms a $D$-region decomposition of $G$ which contradicts
the maximality of $\mathcal{R}$. Let us now consider the other case $v\notin
D$. Since $D$ is a double dominating set of $G$, there exists $x,y\in D$ such
that $vx,vy\in E$. In this case, the path $P=(x,v,y)$ forms a region $R$. Here
also, $\mathcal{R}\cup R$ forms a $D$-region decomposition of $G$ which
contradicts the maximality of $\mathcal{R}$. Thus each vertex of $V$ is in
$V(\mathcal{R})$, that is, $V\subseteq V(\mathcal{R})$. Thus
$V=V(\mathcal{R})$. ∎
It is obvious that, for a plane graph $G=(V,E)$ with a double dominating set
$D$, there exists a maximal $D$-region decomposition $\mathcal{R}$. Based on
Lemma 12, we propose a greedy algorithm to compute a maximal $D$-region
decomposition, which is given in Algorithm 1. The algorithm basically ensures
the properties of the region decomposition mentioned in Definitions 10 and 11.
Input: A plane graph $G=(V,E)$ and a double dominating set $D\subseteq V$.
Output: A maximal $D$-region decomposition $\mathcal{R}$ of $G$.
begin
$V_{used}\leftarrow\emptyset$, $\mathcal{R}\leftarrow\emptyset$;
while _$V_{used}\neq V$_ do
Select a vertex $x$ from $V\setminus V_{used}$;
Consider the set $\mathcal{R}_{x}$ of all regions $S$ with the following
properties:
1. 1.
$S$ is a region between $u$ and $v$, where $u,v\in D$.
2. 2.
$S$ contains $x$.
3. 3.
no vertex from $D\setminus\\{u,v\\}$ is in $V(S)$.
4. 4.
$(S\cup R)\subseteq(\partial(S)\cup\partial(R))$ for all $R\in\mathcal{R}$.
Choose a region $S_{x}\in\mathcal{R}_{x}$ which is maximal in terms of
vertices;
$\mathcal{R}\leftarrow\mathcal{R}\cup\\{S_{x}\\}$;
$V_{used}\leftarrow V_{used}\cup V(S_{x})$;
$\textnormal{{return}}(\mathcal{R})$;
Algorithm 1 REGION$\\_$DECOMPOSITION$(G,D)$
Clearly Algorithm 1 output a maximal $D$-region decomposition in polynomial
time. Next, we show that for a given plane graph $G$ with a double dominating
set $D$, every maximal $D$-region decomposition contains at most $O(|D|)$ many
regions. For that purpose, we observe that a $D$-region decomposition induces
a graph in a very natural way.
###### Definition 13
The induced graph $G_{\mathcal{R}}=(V_{\mathcal{R}},E_{\mathcal{R}})$ of a
$D$-region decomposition $\mathcal{R}$ of $G$ is the graph with possible
multiple edges which is defined as follows: $V_{\mathcal{R}}=D$ and
$E_{\mathcal{R}}=\\{(u,v)|$there is a region $R(u,v)\in\mathcal{R}$ between
$u,v\in D\\}$.
Note that, since by Definition 11 the regions of a $D$-region decomposition do
not intersect, the induced graph $G_{\mathcal{R}}$ of a $D$-region
decomposition $\mathcal{R}$ is a planar graph with multiple edges. Next we
bound the number of regions in a maximal $D$-region decomposition using the
concept of _thin planar graph_ following Alber et al. [AFN04a].
###### Definition 14
A planar graph $G=(V,E)$ with multiple edges is thin if there exists a planar
embedding such that if there are two edges $e_{1},e_{2}$ between a pair of
distinct vertices $v,w\in V$, then there must be two further vertices
$u_{1},u_{2}\in V$ which sit inside the two disjoint areas of the plane that
are enclosed by $e_{1}$ and $e_{2}$.
###### Lemma 15
Let $D$ be a double dominating set of a planar graph $G=(V,E)$. Then the
induced graph $G_{\mathcal{R}}=(V_{\mathcal{R}},E_{\mathcal{R}})$ of a maximal
$D$-region decomposition $\mathcal{R}$ of $G$ is a thin planar graph.
###### Proof.
Let $R_{1}$ and $R_{2}$ be two regions between two vertices $v,w\in D$ and
$e_{1}$ and $e_{2}$ be the corresponding multiple edges between two vertices
$v,w\in V_{\mathcal{R}}$. Let $A$ be an area enclosed by $e_{1}$ and $e_{2}$.
If $A$ contains a vertex $u\in D$, we are done. Suppose there is no vertex of
$D$ in $A$. Now consider the following cases:
There is no vertex from $V\setminus D$ in $A$:
In this case, by combining the regions $R_{1}$ and $R_{2}$, we can form a
bigger region which is a contradiction to the maximality of $\mathcal{R}$.
There is a vertex $x\in(V\setminus D)$ in $A$:
In this case, if $x$ is double dominated by $v$ and $w$, then again we can
combine the two regions $R_{1}$ and $R_{2}$ to get a bigger region. So, assume
that $x$ is dominated by some vertex $u$ other than $v$ and $w$. Since $G$ is
planar, $u$ must be in $A$ which contradicts the fact that $A$ does not
contain any vertex from $D$.
Hence, combining both the cases we see that $G_{\mathcal{R}}$ is a thin planar
graph. ∎
In [AFN04a], it is proved that for a thin planar graph $G=(V,E)$, we have
$|E|\leq 3|V|-6$. Hence we have the following lemma.
###### Lemma 16
For a plane graph $G$ with a double dominating set $D$, every maximal
$D$-region decomposition $\mathcal{R}$ contains at most $3|D|$ many regions.
Now, if we can bound the number of vertices that belongs to any region
$R(u,v)$ of a maximal $D$-region decomposition $\mathcal{R}$ by some constant
factor, we are done. However, achieving this constant factor bound is not
possible for any plane graph $G$. But in a reduced plane graph, we can obtain
this bound, as shown in the following lemma.
###### Lemma 17
A region $R$ of a plane reduced graph contains at most $6$ vertices, that is,
$|V(R)|\leq 6$.
###### Proof.
Let $R$ be the region between $u$ and $v$ and $\partial(R)=\\{u,x,v,y\\}$.
First note that $R$ contains at most two vertices from
$N_{G^{\prime}}^{1}(u,v)$ and the only possibility of such vertices are $x$
and $y$. If there exists a vertex $w\in N_{G^{\prime}}^{1}(u,v)$, apart from
$x$ and $y$, then $w$ has to have a neighbor $z\notin N_{G^{\prime}}(u,v)$.
$z$ should be inside the region $R$ and hence, cannot be double dominated.
Now, because of the reduction rule, we can say that
$\left|N_{G^{\prime}}^{3}(u,v)\right|\leq 1$. We consider the two cases:
Case I ($\left|N_{G^{\prime}}^{3}(u,v)\right|=1$):
In this case, $\left|N_{G^{\prime}}^{2}(u,v)\right|=\emptyset$ by the
reduction rule. Hence, $|V(R)|\leq 5$.
Case II ($\left|N_{G^{\prime}}^{3}(u,v)\right|=0$):
In this case, we claim that there can be at most two vertices from
$N_{G^{\prime}}^{2}(u,v)$. If possible, let $p,q,r\in
N_{G^{\prime}}^{2}(u,v)$. Now all these three vertices must be adjacent to
either $x$ or $y$, which is not possible because of planarity. Hence, in this
case $|V(R)|\leq 6$.
∎
First observe that, for a reduced graph $G^{\prime}$ with a minimum double
dominating set $D$, by Lemma 16, there exists a maximal $D$-region
decomposition $\mathcal{R}$ with at most $3\cdot\gamma_{2}(G^{\prime})$
regions. Also by Lemma 12, we have $V^{\prime}=V(\mathcal{R})$ and by Lemma
17, we have for each region $|V(R)|\leq 6$. Thus we have
$|V^{\prime}|=|V(\mathcal{R})|=|\cup_{R\in\mathcal{R}}V(R)|\leq\sum_{R\in\mathcal{R}}\left|V(R)\right|\leq
6\cdot|\mathcal{R}|\leq 18\cdot\gamma_{2}(G^{\prime})$. Hence, we have the
following theorem.
###### Theorem 18
For a reduced planar graph $G^{\prime}=(V^{\prime},E^{\prime})$, we have
$|V^{\prime}|\leq 18\cdot\gamma_{2}(G^{\prime})$, that is, Double Domination
Problem on planar graph admits a linear kernel.
### 4.2 Liar’s and $k$-tuple domination
We first show that the number of vertices in a plane graph,
$|V|=O(\gamma_{LR}(G))$. In this respect, first we note that both the results
in Lemma 12 and Lemma 16 are valid for any plane graph $G$ and any double
dominating set $D$. Since every liar’s dominating set is also a double
dominating set, similar type of results hold for any plane graph $G$ and any
liar’s dominating set $L$. We claim that the number of vertices in a region
$R$ of a $L$-region decomposition is bounded above by a constant. Let $R$ be a
region between $u$ and $v$ and $\partial(R)=\\{u,x,v,y\\}$. Note that in
$V(R)$ there are two vertices ($u$ and $v$) from $L$. Now, if there exists two
vertices $p,q\in V(R)\setminus\partial(R)$, then for the pair $p$ and $q$
condition (ii) of liar’s domination is violated. Hence, there is at most one
vertex in $V(R)\setminus\partial(R)$. Therefore, $|V(R)|\leq 5$. Thus we have
$|V|=|V(\mathcal{R})|=|\cup_{R\in\mathcal{R}}V(R)|\leq\sum_{R\in\mathcal{R}}|V(R)|\leq
5\cdot|\mathcal{R}|\leq 15\cdot|L|\leq 15\cdot\gamma_{LR}(G).$ Hence, we have
the following theorem
###### Theorem 19
For a planar graph $G=(V,E)$, $|V|\leq 15\cdot\gamma_{LR}(G)$.
Since every $k$-tuple dominating set for $k\geq 3$ is a liar’s dominating set,
we can use Theorem 19. But, we can improve the constant a little bit.
###### Theorem 20
For a planar graph $G=(V,E)$, $|V|\leq 12\cdot\gamma_{k}(G)$, where $k\geq 3$.
###### Proof.
Let $D$ be a minimum $k$-tuple dominating set of $G=(V,E)$. Since every
$k$-tuple dominating set is a double dominating set, by Lemma 16 we can form a
maximal $D$-region decomposition $\mathcal{R}$ of $G$ containing at most
$3\cdot|D|$ many regions. Again by Lemma 12, we have $V=V(\mathcal{R})$. Since
each region contains only two vertices of $D$, we have $|V(R)|\leq 4$.
Otherwise there exists one vertex in $V(R)$ which is not dominated by $k$
vertices of $D$. Hence $|V|\leq 4\cdot|\mathcal{R}|\leq 12\cdot|D|\leq
12\cdot\gamma_{k}(G)$. ∎
## 5 Linear kernels for bounded genus graphs
In this section, we extend our results to bounded genus graphs to show that
$k$-Tuple Domination Problem and Liar’s Domination Problem admit a linear
kernel. The notations in this section follow Section 2.1.
For double domination problem, we apply the same reduction rule on a graph $G$
with bounded genus $g$ to obtain the reduced graph $G^{\prime}$. Note that the
reduced graph $G^{\prime}$ is also of bounded genus $g$. Let $G=(V,E)$ be an
$n$-vertex $\Sigma$-embedded graph. It is easy to observe that, since
$\sum_{v\in V}deg_{G}(v)=O(n+\textup{eg}(\Sigma))$, the reduced graph
$G^{\prime}=(V^{\prime},E^{\prime})$ can be computed in
$O(n^{3}+n^{2}\cdot\textup{eg}(\Sigma))$ time, where $|V|=n$. Next we show
that $|V^{\prime}|=O(\gamma_{2}(G^{\prime})+g)$ which implies Double
Domination Problem admits a linear kernel in bounded genus graphs.
To prove the above, we consider two cases. In the first case, we assume that
the reduced $\Sigma$-embedded graph has representativity strictly greater than
$4$. In the case when $rep(G)\leq 4$, we go by induction on the Euler genus of
surface $\Sigma$. In the first case, the graphs are locally planar, i.e., all
the contractable noose are of length less or equal to $4$. Since the boundary
of the regions in planar case is less than or equal to $4$, the boundary
$\partial(R)$ of any region $R$ of a $D$-region decomposition $\mathcal{R}$ is
contractible. Hence the proof in the planar case can be extended in this case.
Hence we have the following lemma.
###### Lemma 21
Let $G^{\prime}=(V^{\prime},E^{\prime})$ be a reduced $\Sigma$-embedded graph
where $\textup{rep}(G^{\prime})>4$. Then $|V^{\prime}|\leq
18(\gamma_{2}(G^{\prime})+\textup{eg}(\Sigma))$.
###### Proof.
Let $D$ be a double dominating set of $G^{\prime}$ and $\mathcal{R}$ is a
maximal $D$-region decomposition of $G^{\prime}$. Forming a induced graph,
$G_{\mathcal{R}}$ as in case of double domination problem in planar graphs
(Section 4.1.2), we have $|\mathcal{R}|\leq 3\cdot(|D|+eg(\Sigma))$. Also, in
this case, every vertex of $V^{\prime}$ belongs to at least one region of
$\mathcal{R}$ and for a region $R$, $|V(R)|\leq 6$. Hence, we have
$|V^{\prime}|\leq 18(\gamma_{2}(G^{\prime})+\textup{eg}(\Sigma))$. ∎
Next consider the case where $3\leq\textup{rep}(G^{\prime})\leq 4$. For a
noose $N$ in $\Sigma$, we define the graph $G_{N}=(V_{N},E_{N})$ as follows.
First we consider the graph $\mathcal{G}$ obtained from
$G^{\prime}=(V^{\prime},E^{\prime})$ by cutting along $N$. Then for every
$v\in N\cap V^{\prime}$ if $v^{i}$, $i=1,2$, is not adjacent to a pendant
vertex, then we add a pendant vertex $u^{i}$ adjacent to $v^{i}$ to form
$G_{N}$. Clearly $G_{N}$ has genus less than that of $G^{\prime}$. If we add
all the vertices of $V_{N}\setminus V^{\prime}$ to a double dominating set $D$
of $G^{\prime}$, then we clearly obtain a double dominating set of $G_{N}$ and
as, $\textup{rep}(G^{\prime})\leq 4$, $|V_{N}\setminus V^{\prime}|\leq 16$.
Hence, $\gamma_{2}(G_{N})\leq\gamma_{2}(G^{\prime})+|N\cap
V^{\prime}|\leq\gamma_{2}(G^{\prime})+16$. Also note that if $G^{\prime}$ is a
reduced graph, then so is $G_{N}$. Using these facts, we prove that Double
Domination Problem possesses a linear kernel when restricted to graphs with
bounded genus.
###### Lemma 22
For any reduced $\Sigma$-embedded graph $G^{\prime}=(V^{\prime},E^{\prime})$
with eg$(\Sigma)\geq 1$, $|V^{\prime}|\leq
18(\gamma_{2}(G^{\prime})+32\cdot\textup{eg}(\Sigma)-16)$.
###### Proof.
We prove this result by induction on $\textup{eg}(\Sigma)$. Suppose
$\textup{eg}(\Sigma)=1$. If $\textup{rep}(G^{\prime})>4$, then the result
follows from Lemma 21. Otherwise Lemma 1 implies that the graph $G_{N}$,
described above, is planar. Hence by Theorem 18, we have $|V_{N}|\leq
18\cdot\gamma_{2}(G_{N})$. Thus $|V^{\prime}|\leq|V_{N}|\leq
18(\gamma_{2}(G^{\prime})+16)$.
Assume that $|V^{\prime}|\leq
18(\gamma_{2}(G^{\prime})+32\cdot\textup{eg}(\Sigma)-16)$ for any
$\Sigma$-embedded reduced graph $G^{\prime}$ with eg$(\Sigma)\leq g-1$.
Consider a reduced $\Sigma$-embedded graph $G^{\prime}$ with eg$(\Sigma)=g$.
Now if rep$(G^{\prime})>4$, then again by Lemma 21, we are done. Hence assume
that rep$(G^{\prime})\leq 4$. By Lemma 1, either $G_{N}$ is the disjoint union
of graphs $G_{1}$ and $G_{2}$ that can be embedded in surfaces $\Sigma_{1}$
and $\Sigma_{2}$ such that eg$(\Sigma)=$ eg$(\Sigma_{1})+$ eg$(\Sigma_{2})$
and eg$(\Sigma_{i})>0$, $i=1,2$ (this is the case when $N$ is surface
separating curve), or $G_{N}$ can be embedded in a surface with Euler genus
strictly smaller than eg$(\Sigma)$ (this holds when N is not surface
separating).
Let us consider the case where $G_{N}$ is the disjoint union of graphs
$G_{1}=(V_{1},E_{1})$ and $G_{2}=(V_{2},E_{2})$ that can be embedded in
surfaces $\Sigma_{1}$ and $\Sigma_{2}$. Since eg$(\Sigma_{i})\leq g-1$ for
$i=1,2$, we can apply the induction hypothesis on $G_{i}$. Thus we have,
$\displaystyle|V^{\prime}|\leq|V_{N}|$ $\displaystyle=|V_{1}|+|V_{2}|$
$\displaystyle\leq\sum_{i=1}^{2}18(\gamma_{2}(G_{i})+32\cdot\textup{eg}(\Sigma_{i})-16)$
$\displaystyle\leq 18(\gamma_{2}(G_{N})+32\cdot\textup{eg}(\Sigma)-32)$ as
$G_{1}$ and $G_{2}$ are disjoint $\displaystyle\leq
18(\gamma_{2}(G^{\prime})+32\cdot\textup{eg}(\Sigma)-16).$
Next we consider the case where $G_{N}$ can be embedded in a surface
$\Sigma^{\prime}$ with Euler genus strictly smaller than $g$. In this case
too, we can apply induction hypothesis on $G_{N}$. Thus we have,
$\displaystyle|V^{\prime}|\leq|V_{N}|$ $\displaystyle\leq
18(\gamma_{2}(G_{N})+32\cdot\textup{eg}(\Sigma^{\prime})-16)$
$\displaystyle\leq 18(\gamma_{2}(G_{N})+32\cdot(\textup{eg}(\Sigma)-1)-16)$
$\displaystyle\leq 18(\gamma_{2}(G^{\prime})+32\cdot\textup{eg}(\Sigma)-32)$
as $\gamma_{2}(G_{N})\leq\gamma_{2}(G^{\prime})+16$ $\displaystyle\leq
18(\gamma_{2}(G^{\prime})+32\cdot\textup{eg}(\Sigma)-16).$
Thus we have proved that, $|V^{\prime}|\leq
18(\gamma_{2}(G^{\prime})+32\cdot\textup{eg}(\Sigma)-16)$ for every
$\Sigma$-embedded graph $G^{\prime}=(V^{\prime},E^{\prime})$. ∎
Hence by Lemma 21 and Lemma 22, we have the main result of this subsection.
###### Theorem 23
Double Domination Problem admits a linear kernel for bounded genus graphs.
For liar’s domination problem, by Theorem 19, we have $|V|\leq
15\cdot\gamma_{LR}(G)$ in case of a planar graph $G=(V,E)$. Proceeding exactly
in the same way as in the case of double domination, we can have the following
theorem for a $\Sigma$-embedded graph $G$.
###### Theorem 24
Let $G=(V,E)$ be a $\Sigma$-embedded graph. Then $|V|\leq
15(\gamma_{LR}(G)+32\cdot\textup{eg}(\Sigma))$.
Since for any graph that admits a $k$-tuple dominating set ($k\geq 3$),
$\gamma_{LR}(G)\leq\gamma_{k}(G)$, we have the following corollary of Theorem
24.
###### Corollary 25
For a $\Sigma$-embedded graph $G=(V,E)$, $|V|\leq
15(\gamma_{k}(G)+32\cdot\textup{eg}(\Sigma))$.
## 6 Conclusion
In this paper, we first have proved that $k$-Tuple Domination Problem and
Liar’s Domination Problem are W[2]-hard for general graphs. Then we have shown
that these two problems admit linear kernel for planar graphs and also for
bounded genus graphs. It would be interesting to look for other graph classes
where these problems admit efficient parameterized algorithms.
## Acknowledgements
The authors want to thank Venkatesh Raman and Saket Saurabh for some of their
nice introductory expositions to parametrization.
## 7 Appendix
###### Lemma 26
Liar’s Domination Problem is NP-complete for planar graphs.
###### Proof.
The reduction is from Domination Problem in planar graphs, which is known to
be NP-complete[GJ79]. Let $G=(V,E)$ be a planar graph with
$V=\\{v_{1},v_{2},\ldots,v_{n}\\}$ and $k$ be an integer. We construct an
instance $G^{\prime}=(V^{\prime},E^{\prime})$ and $k^{\prime}$ of Liar’s
Domination Problem as follows: We add a set of $3n$ new vertices
$S=\\{x_{i},y_{i},z_{i}|1\leq i\leq n\\}$ to the vertex set of $V$, i.e.,
$V^{\prime}=V\cup S$ and the edge set of $G^{\prime}$ is given by
$E^{\prime}=E\cup\\{v_{i}x_{i},x_{i}y_{i},y_{i}z_{i}|1\leq i\leq n\\}$. Note
that, since $G$ is planar, so is $G^{\prime}$. Also assume that
$k^{\prime}=k+3n$. In [Sla09], it is proved that $G$ has a dominating set of
cardinality at most $k$ if and only if $G^{\prime}$ has a liar’s dominating
set of cardinality at most $k^{\prime}=k+3n$.
Thus, Liar’s Domination Problem is NP-complete for planar graphs. ∎
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|
arxiv-papers
| 2013-09-21T11:07:24 |
2024-09-04T02:49:51.268082
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Arijit Bishnu, Arijit Ghosh and Subhabrata Paul",
"submitter": "Arijit Ghosh",
"url": "https://arxiv.org/abs/1309.5461"
}
|
1309.5468
|
11institutetext: Variable Energy Cyclotron Centre,1/AF Bidhannagar, Kolkata
700064, India
S.N. Bose National Centre for Basic Sciences, Block JD, Salt Lake, Kolkata
700098, India
Time series analysis in nonlinear dynamics Scaling phenomena in complex
systems Levy flights
# Universal Scaling Property of System Approaching Equilibrium
P. Barat 11 A. Giri 11 M. Bhattacharya 11 Nilangshu K. Das 11 A. Dutta 221122
###### Abstract
In this Letter we show that the diffusion kinetics of kinetic energy among the
atoms in non-equilibrium crystalline systems follows universal scaling
relation and obey Levy-walk properties. This scaling relation is found to be
valid for systems no matter how far they are driven out of equilibrium.
###### pacs:
05.45.Tp
###### pacs:
89.75.Da
###### pacs:
05.40.Fb
Non-equilibrium systems are ubiquitous in nature. Equilibrium systems are
ideal and can only be achieved in the laboratory. Attempts have been made to
understand the dynamics of non-equilibrium systems by linear response theory
[1] albeit its domains of validity are restricted to the linear response
regime. There exists no general formalism to deal with systems that are far
from equilibrium. For a system out of equilibrium, the probability of a given
microstate evolves continuously with time. In the long time limit the system
reaches a stationary state in which the probability measure over the
configuration space converges to a constant distribution. Non-equilibrium
systems display fluctuations which are less sensitive to the conditions of the
surroundings and carry information about the dynamics of its present state. In
the last decade certain general relations have been discovered which are valid
for non-equilibrium systems and are independent of how far the system is
driven out of equilibrium. These results include the Jarzynski equality [2, 3]
and the fluctuation theorems [4, 5, 6, 7, 8, 9, 10, 11, 12]. They have been
verified for a variety of systems theoretically as well as experimentally [13,
14, 15, 16, 17]. After the work by Crooks [7] and Seifert [8], it is now
understood that many of these relations are closely related to the path
probability of the system’s trajectory. In the absence of any general theory
to understand the dynamics of a system far from equilibrium, one approach will
be to take a simple but nontrivial model system and try to understand its
dynamics when it approaches equilibrium from its non-equilibrium state. The
general dynamical behavior of a non-equilibrium system will consist of
superposition of various dynamics on well speared time scales which compel
several thermodynamic parameters of the system like heat, work, internal
energy, kinetic energy of the particles of the system to fluctuate. To
understand the general features of the dynamics one has to apply statistical
analysis to these fluctuations. In this Letter an attempt has been made to
study a universal scaling relation for systems driven far from equilibrium by
analyzing its fluctuation properties. The non-equilibrium systems of study are
generated by molecular dynamics simulation technique.
Dynamical systems arising from diversified disciplines of science can be
quantified in a unified way from their scale invariance properties. Scafetta
$et$ $al.$ [18] introduced two complementary scaling analysis methods:
Diffusion Entropy Analysis (DEA) and Finite Variance Scaling Method (FVSM) to
evaluate correct scaling that prevails in complex dynamical systems. DEA
studies the scaling exponent $\delta$ of the probability distribution instead
of its moment, based on the evaluation of Shannon entropy $S(t)$ of a time
series $\\{\xi_{i}\\}$. The DEA has been successfully applied to time series
of different kinds of dynamical systems [19]. In DEA analysis, the numbers in
the time series, $\\{\xi_{i}\\}$, are thought to be the fluctuations of a
diffusion trajectory [19] arising from the probability density function (PDF),
$p(x,t)$, of the corresponding diffusion process. Here $x$ denotes the
variable collecting the fluctuations and is referred as the diffusion
variable. The scaling property of $p(x,t)$ is evaluated by means of the sub
trajectories $x_{n}(t)=\sum_{i=0}^{t}\xi_{i+n}$ with n = 0,1,….etc. If the
scaling condition $p(x,t)=t^{-\delta}F(xt^{-\delta})$ holds true, it is easy
to prove that the entropy $S(t)$ is given by
$S(t)=-\int^{+\infty}_{-\infty}p(x,t)ln[p(x,t)]dx=A+\delta ln(t)$, where $A$
is a constant. This indicates that in the case of a diffusion process with a
scaling PDF, its entropy $S(t)$ increases linearly with $ln(t)$. One can also
examine the scaling properties of the second moment of the same process by
FVSM. One version of FVSM is the Standard Deviation Analysis (SDA) [19], which
is based on the evaluation of the standard deviation $D(t)$ of the variable
$x$, and yields $D(t)=\left[\langle x^{2};t\rangle-\langle
x;t\rangle^{2}\right]^{\frac{1}{2}}\propto t^{H}$ [19]. The exponent $H$ is
interpreted as the scaling exponent. For random noise with finite variance,
the diffusion distribution $p(x,t)$ will converge, according to the central
limit theorem, to a Gaussian distribution with $H=\delta=0.5$. If
$H\neq\delta$, the scaling represents anomalous behavior. Levy-walk is a kind
of anomalus diffusion which is obtained by generalizing the central limit
theorem [20]. In this particular kind of diffusion process the scaling
exponents $H$ and $\delta$ are found to obey the relation $\delta=(3-2H)^{-1}$
instead of being equal [19].
Table 1: The parameters used for the MD simulations. Material | Si | Ge | Cu | Solid Ar | Mo | Fe | Al |
---|---|---|---|---|---|---|---|---
Simulation cell size (in unit cell) | $25^{3}$ | $15^{3}$ | $20^{3}$ | $20^{3}$ | $25^{3}$ | $25^{3}$ | $20^{3}$ |
Number of atoms in simulation cell | 125000 | 27000 | 32000 | 32000 | 31250 | 31250 | 32000 |
Simulation time step (fs) | 0.25 | 0.5 | 0.5 | 2.0 | 0.5 | 0.5 | 0.5 |
Nature of inter-atomic potential used | SW111SW (Stillinger-Weber) | SW | EAM222 EAM (Embedded Atom Model) | LJ333LJ (Lennard-Jones) | FS444FS (Finnis-Sinclair) | FS | GLUE |
Total number of data taken | 6000 | 3000 | 3000 | 10000 | 3000 | 3000 | 3000 |
Radius ‘r’ (Å) | 8.8250 | 9.0809 | 5.4766 | 7.9952 | 6.2944 | 5.7330 | 6.1566 |
No of atoms in the sphere of radius ‘r’ | 150 | 147 | 55 | 55 | 63 | 63 | 55 |
To generate the non-equilibrium model system we have utilized the molecular
dynamics (MD) simulation technique in an innovative way. Depending on the
interaction potential model, a typical MD simulation computes the trajectories
of atoms in a system by solving Newton’s equations of motion numerically. In
our study, the simulations have been performed for elements having different
crystal structures like Silicon (Si) and Germanium (Ge) with diamond cubic
structure, Iron (Fe) and Molybdenum (Mo) with body centered cubic structure
and Copper (Cu), Aluminium (Al) and solid Argon (Ar) having face centered
cubic structure. Thus, a variety of interaction potential models have been
employed (refer Table 1) in the simulations. For all these elements, periodic
boundary conditions are imposed along three directions (x,y,z) in the
simulation cell. Initially, all the systems are equilibrated using constant
energy and volume ensemble (NVE) at 100 K for 2.5 ps except for Ar, where the
equilibration run is performed at 30K for 10 ps. The kinetic energy in the
equilibrium state will follow Boltzmann distribution. The period of time for
equilibration is chosen to be sufficient to bring the total energy of the
system to divide approximately equally between kinetic energy (KE) and
potential energy (PE). Details of the parameters used in the simulations like
simulation cell size, number of atoms taken in the simulation cell, the time
steps after which simulation data are recorded and the nature of the inter-
atomic potentials used for the elements under study are given in Table I. To
generate the non-equilibrium state, the three components of velocities of
individual atom in the simulation cell were changed to random values such that
the instantaneous KE of the atom increases however the average temperature of
the system corresponds to T = 500 K (T = 70 K for Ar). The atoms in the
simulation cell are then allowed to equilibrate. After sufficiently long
period of time, the system is observed to attain equilibrium temperature of
300 K (50 K for Ar) and the average kinetic energy becomes equal to the
average potential energy of the respective systems. The variations of
temperature (equivalent to KE) with time of the systems of Ar and Cu leading
to equilibration are shown in Fig. 1. Atomistic simulation data are recorded
for each time step for the entire equilibration time to study the dynamic
nature of the equilibration process.
Figure 1: Variations of temperature (equivalent to KE ) with time of the
systems of (a) Ar and (b) Cu leading to the equilibration
The molecular dynamics simulations are carried out using MD++ simulation
package [21].
Few atoms in the simulation cell are identified in a sphere of radius ‘r’
whose centre is chosen to be almost at the middle of the cell, far away from
the surface. The values of ‘r’ and the number of atoms in the sphere for each
simulation are given in Table 1. The KE in the consecutive time steps for the
atoms in the sphere are used for further analysis. In these simulations the
atoms in the simulation cell are kept isolated from the environment and hence
there was no dissipation of energy, only there is a redistribution of KE and
PE among the atoms. DEA and SDA analyses [22] were performed for the time
series data of the KE for an individual atom. Typical plots of DEA and SDA
analysis from which $\delta$ and $H$ are calculated for a Copper atom are
shown in Fig. 2. The slopes of the curves evaluate $\delta$ and $H$. DEA and
SDA analyses were performed for ten atoms randomly chosen from the atoms
confined in the sphere of radius ‘r’ and the average of values of $\delta$ and
$H$ are given in Table 2. Simulation cell sizes, number of atoms in the cell,
sampling time, the crystal structure, and the nature of the interaction
potential were varied to find the exact nature of the scaling.
Table 2: Mean values of the scaling exponents $\delta$ and H obtained from ten atoms chosen randomly in the shell of radius ‘r’. Exponents $\delta$ and $H$ are evaluated from the time series of the evolution of KE for a single atom. Material | Non-equilibrium | $\delta$ | $H$ | $\left[\left(\delta-\frac{1}{3-2H}\right)/\delta\right]\times 100$ |
---|---|---|---|---|---
| Temperature(K) | | | |
Si | 500 | 0.953±0.016 | 0.968±0.008 | 1.380 |
Si | 800 | 0.898±0.025 | 0.901±0.025 | 7.046 |
Si ( single atom ) | 500 | 0.943 | 0.970 | 0.042 |
Ge | 500 | 0.930±0.021 | 0.953±0.026 | 1.712 |
Ge | 800 | 0.930±0.020 | 0.960±0.015 | 0.438 |
Ar | 70 | 0.968±0.010 | 0.968±0.005 | 2.908 |
Ar ( single atom ) | 70 | 0.963 | 0.972 | 1.665 |
Cu | 500 | 0.920±0.020 | 0.955±0.022 | 0.279 |
Cu ( single atom ) | 500 | 0.935 | 0.978 | 2.444 |
Fe | 500 | 0.925±0.018 | 0.964±0.022 | 0.847 |
Fe ( single atom ) | 500 | 0.959 | 0.974 | 0.879 |
Mo | 500 | 0.935±0.016 | 0.964±0.021 | 0.231 |
Al | 500 | 0.935±0.008 | 0.966±0.013 | 0.142 |
Figure 2: (a) DEA and (b) SDA of the variation of kinetic energy against time
data obtained from a Copper atom in an ensemble when brought to a non-
equilibrium state at 500 K from 100 K and allowed to equilibrate at 300 K
To understand the process of equilibration when a single excited atom
interacts with an ensemble of atoms that are in equilibrium, another kind of
exercise was performed by MD simulation. In this case from the equilibrated
ensembles of Si, Cu, Fe at 300 K and solid Ar at 50 K one atom in the middle
of the cell was excited to a KE corresponding to 500 K and 70 K respectively.
The excited atom is then allowed to equilibrate and the KE possesses by the
atom at different time steps was monitored. In these simulations the energy of
the simulation cells was also a constant of motion. All the equilibrated atoms
in the cell form a heat bath and the excited atom equilibrate by interacting
with this heat bath by sharing its excess KE. The time series thus obtained of
the KE of the excited atom between consecutive time steps was used to
understand the diffusion process of KE and the nature of the dynamical process
responsible to bring the atom to the equilibrium state. The values of $\delta$
and $H$ obtained by DEA and SDA analysis from the time series are given in
Table 2.
Atoms in the simulation cell, when brought to a non-equilibrium state by
enhancing their KE, try to redistribute their excess KE by the process of
diffusion. This process of redistribution of KE among the atoms changes their
mean positions of vibration and consequently their PE. Here all the atoms in
the cell participate together in the diffusion process. Thus by analyzing the
variation of the KE with time for a single atom in the cell one can account
for the modality of this diffusion process and the nature of the
equilibration. The change in the KE for an atom in the $i^{th}$ time step is
given by $\Delta E=\left(\frac{2E_{i}}{m}\right)^{\frac{1}{2}}F_{i}\Delta t$
where $F_{i}$ is force on the atom at the $i^{th}$ time step. The force field
experience by an atom in the cell will consist of three parts. The
deterministic force, arising from the nearest neighbor interaction potential,
and it plays the role of the external force $F(t)$ acting on the atom. There
will be energy exchange between the atom and the surrounding atoms in a result
of which the atom loses a part of its KE for exciting various degrees of
freedom of the atoms in the ensemble as well there will be increase in the PE
in the cost of KE of the atom. This can be described with help of a frictional
force $F_{f}(t)$ acting on the atom. Besides loosing KE due to frictional
effect there is a possibility of gaining KE in the form of random movement of
the atom due to interactions from all other atoms in the cell and can be
modeled by a random force $\Gamma(t)$ acting on the atom. $\Gamma(t)$ will
have the property as $\langle\Gamma(t)\rangle=0$ and
$\langle\Gamma(t_{1})\Gamma(t_{2})\rangle=Ag(\Delta t)$ where $\Delta
t=(t_{2}-t_{1})$, the function $g(t)$ dies down rapidly with $t$ and $A$ is
the strength of the random force and is a function of temperature T. Thus the
dynamics of the atoms will follow Langevin type equation [23]. To see the
effect of this random force field on the scaling behavior, the ensembles of Si
and Ge atoms were exited to a higher temperature i.e. at 800 K and allowed to
equilibrate at 450 K. The $\delta$ and $H$ values for the time series of the
KE for these cases are given in Table 2. As the systems were allowed to
equilibrate of its own without any defined protocol, forces acting on the
atoms are time dependent through the dynamical process of equilibration. The
force field experienced by an atom in the
Figure 3: Probability distributions of the fluctuating forced fields
experienced by (a) an Argon and (b) a Silicon atom when brought to a non-
equilibrium state at 500 K from 100 K and allowed to equilibrate at 300 K
simulation cell in the $i^{th}$ time step is calculated from its velocities,
at $i-1$, $i$, $i+1^{th}$ time steps. Typical mean values of these forces are
$1.011654\times 10^{-10}$ N and $24.7088\times 10^{-10}$ N for solid Argon and
Si respectively. The fluctuating part of the force field at different time
steps was obtained by subtracting the mean values. Typical probability
distributions for the fluctuating part of the force field in case of solid Ar
and Si atoms are shown in Fig. 3. The distributions are Gaussian with center
$-6.69\times 10^{-12}$ N and width $9.59\times 10^{-11}$ N for Ar and the
corresponding values for Si are $-4.98\times 10^{-11}$ N and $1.95\times
10^{-9}$ N respectively. To understand the correlation of the fluctuating
force field, vectors $\\{\lvert X_{i}\rangle\\}$ are generated from the time
series data of the fluctuating force field by taking consecutive n time steps
for constructing each vector. The eigen value spectrum of the covariance
matrix $\Sigma_{x}$ of these vectors are calculated to establish the nature of
the correlation. Fig. 4 shows typical eigen value spectra for Ar and Si cases.
The flatness of the spectra confirms that the fluctuating force fields are
uncorrelated.
Figure 4: Eigen value spectra of the covariance matrices constructed from the
time series of the fluctuating forced fields experienced by (a) an Argon and
(b) a Silicon atom when brought to a non-equilibrium state at 500 K from 100 K
and allowed to equilibrate at 300 K
The motion of the atoms in the simulation cell is thus governed by these
forces and dictates the variation of KE of the atoms. The memory effect of the
initial velocity will die down with time. In the long time limit the system
equilibrates leading to equipartition of energy. The characteristic of the
variation of KE in terms of time series indicates the evolution of the non-
equilibrium state. To understand the universality of this evolution for
different dynamical processes as generated by MD simulations, the Shannon
entropy of the diffusion process of KE of an atom in the cell was obtained by
calculating the PDF from the sub trajectories of this time series. The high
values of $\delta$ and $H$ as given in Table 2 signifies a strong persistence
in the fluctuations of the KE of the atoms. The values of $H$ are always
larger than that of $\delta$ for all cases studied and are seen to fulfill the
Levy-walk diffusion relation within the error bar as shown in Table 2.
The exact inherent dynamics of the process of transmitting excess KE of an
atom to its surroundings could not be revealed from this analysis. However,
this analysis reflected the underlying generic features and physical
principles that are independent of the detailed dynamics or characteristics of
particular model. The diffusion of KE of an atom to its surroundings is a
continuous stationary stochastic process as the probability of the diffusion
trajectories follows a scaling relation. As the exponent $H$ is greater than
0.5 for all cases, the diffusion is anomalous super diffusion. Any diffusion
is a kind of random walk and the Levy-walk is a mathematical model to describe
anomalous super diffusion where the scaling exponent of variance against time
is greater than one. Levy-walks have coupled space-time probability
distributions and are characterized by a cluster of smaller variations of the
random variable with a few large variations between them. This pattern repeats
for all scales. As the process of equilibration of KE of an atom in the
ensemble from its non-equilibrium state has to be very rapid, it cannot be
Brownian type and it should be Levy-walk type as Levy-walk will outperform
Brownian walk during the process of equilibration. The Levy-walk type of
diffusion of KE among the atoms may be due to the following reason. Excess KE
of an atom should disburse locally however there is a finite probability to
transport its excess KE to a distant atom and consequently this atom again
redistributes its KE in the same process. This process is much faster than
normal diffusion when the mean squared value of fluctuation depends on t only.
Mostly three dimensional systems show normal diffusion except in glassy
systems [24]. The anomalous or superdiffusion and Levy-walks have been
observed in various real-life phenomena like fluid flow in rotating annulus
[25], low dimensional heat transport [26], light scattering in porous media
[27] etc. However our findings show that the diffusion of KE for a single atom
in its non-equilibrium state, when embedded in an environment of atoms that
are either in equilibrium or in non-equilibrium state, show superdiffusion and
Levy-walk properties throughout the process of equilibration. In the case when
a single atom interacts with the ensemble of atoms that are in equilibrium we
can assume that the atom is interacting with a thermal bath by absorbing or
releasing KE without appreciable change in the bath state. In this case also
we see the diffusion of KE from the atom is a Levy-walk process.
Atoms of crystals when brought to a non-equilibrium state, try to redistribute
its energy amongst them to converge to an equilibrium distribution. This
fundamental process is governed by the increase in entropy of the trajectories
of the individual atom and should be identical in nature for all types of
elemental atoms and will be independent of the nature of interactions between
the atoms and the perturbation which brought them to non-equilibrium state.
The calculation of Shannon entropy of the diffusion process was made for one
atom. However, the concept of entropy in statistical mechanics is for an
ensemble. But the entropy production for a single trajectory has been
addressed in the literature [8] concerning fluctuation theorems.
Atoms in a lattice possessing KE different from its equilibrium values at any
instant of time try to transport or accept from the nearest neighbors and
arrive at a new value. The time series of this KE were translated in to a
diffusion process in the form of diffusion trajectories and PDF of these
trajectories were estimated. The estimated Shannon entropy production of this
dynamical process with time is found to obey a universal scaling relation.
This scaling relation is exact and valid for systems no matter how far they
are driven out of equilibrium and is independent of the strength of
perturbation that brought the system out of equilibrium.
## References
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* [2] C. Jarzynski, Annu. Rev. Condens. Matter Phys. 2 (2011) 329 .
* [3] C. Jarzynski, Phys. Rev. Lett. 78 (1997) 2690 .
* [4] D. J. Evans, E G. D. Cohen, and G. P. Morriss, Phys. Rev. Lett. 71 (1993) 2401 .
* [5] D. J. Evans and D. J. Searles, Phys. Rev. E 50 (1994) 1645 .
* [6] G. Gallavotti and E. G. D. Cohen, Phys. Rev. Lett. 74 (1995) 2694 .
* [7] G. E. Crooks, Phys. Rev. E 60 (1999) 2721 .
* [8] Udo Seifert, Phys. Rev. Lett. 95 (2005) 040602 .
* [9] J. Kurchan, J. Phys. A: Math. Gen. 31 (1998) 3719 .
* [10] J. L. Lebowitz and H. Spohn, J. Stat. Phys. 95 (1999) 333 .
* [11] C. Maes, Sem. Poincare 2 (2003) 29 .
* [12] P. Gaspard, J. Chem. Phys. 120 (2004) 8898 .
* [13] J. Liphardt, S. Dumont, S.B.Smith, I. Tinoco Jr, and C. Bustamante, Science 296 (2002) 1832 .
* [14] G. M. Wang, E.M.Sevick, Emil Mittag, Debra J. Searles and Denis J. Evans, Phys. Rev. Lett. 89 (2002) 050601 .
* [15] K. Feitosa and N. Menon, Phys. Rev. Lett. 92 (2004) 164301.
* [16] W. I. Goldburg, Y. Y. Goldschmidt and H. Kellay, Phys. Rev. Lett. 87 (2001) 245502.
* [17] N. Garnier and S. Ciliberto, Phys. Rev E 71 (2005) 060101(R).
* [18] N. Scafetta. P. Hamilton and P. Grigolini, Fractals 9 (2001) 193.
* [19] N. Scafetta and P. Grigolini, Phys. Rev. E 66 (2002) 036130, N. Scafetta and B.J.West, Phys. Rev. Lett. 92 (2004) 138501\.
* [20] B.V. Gnedenko, A.N. Kolmogorov, Limit Distributions for Sums of Independent Random Variables (Addison-Wesley, Reading, MA) 1954.
* [21] http//micro.stanford.edu
* [22] The programme for calculating DEA and SDA for the time series data was initially tested for synthetic time series like Gaussian White Noise generated by Box-Muller algorithm, Fractional Gaussian Noise generated based on the reference, Methodology and Computing in Applied Probability 1 (1999) 445 and Levy Noise generated based on the reference Phys. Rev. E 49 (1994) 4677.
* [23] G.W. Ford, M. Kac and P. Mazur, J. Math. Phys. 6 (1965) 504.
* [24] L. Cipelletti and L. Ramos, J. Phys. : Condens. Matter 17 (2005) R253 and the references therein.
* [25] T. H. Solomon, E. R. Weeks and H. L. Swinney, Phys. Rev. Lett. 71 (1993) 3975.
* [26] A. Dhar, K.Saito and B. Derrida, Phys. Rev. E 87 (2013) 010103 (R) and the references therein.
* [27] P. Barthelemy, J. Bertolotti and D. S. Wiersma, Nature 453 (2008) 495.
|
arxiv-papers
| 2013-09-21T12:27:48 |
2024-09-04T02:49:51.278288
|
{
"license": "Public Domain",
"authors": "P. Barat, A. Giri, M. Bhattacharya, Nilangshu K. Das and A. Dutta",
"submitter": "Amal Giri",
"url": "https://arxiv.org/abs/1309.5468"
}
|
1309.5507
|
# Microgrid - The microthreaded many-core architecture
Irfan Uddin
University of Amsterdam, The Netherlands
[email protected]
###### Abstract
Traditional processors use the _von Neumann_ execution model, some other
processors in the past have used the _dataflow_ execution model. A combination
of _von Neuman_ model and _dataflow_ model is also tried in the past and the
resultant model is referred as _hybrid dataflow_ execution model. We describe
a hybrid dataflow model known as the microthreading. It provides constructs
for creation, synchronization and communication between threads in an
intermediate language. The microthreading model is an abstract programming and
machine model for many-core architecture. A particular instance of this model
is named as the microthreaded architecture or the Microgrid. This architecture
implements all the concurrency constructs of the microthreading model in the
hardware with the management of these constructs in the hardware.
###### Contents
1. 1 Introduction
2. 2 Latency tolerance
3. 3 The microthreading model
4. 4 The Microgrid
5. 5 Programming the Microgrid
6. 6 The microthreading model in the context of the Microgrid
7. 7 I/O in the Microgrid
8. 8 Conclusion
## 1 Introduction
Traditional processors are based on the _von Neumann_ execution model. In this
model a sequence of instructions is executed one-by-one and the state of the
program is identified by a single program counter. However, performance can
not be improved by executing instructions sequentially. In order to improve
the performance, _dataflow_ scheduling is used where instructions are executed
based on the availability of data. Theoretically, dataflow models are parallel
execution models, because instructions can be scheduled only with the
availability of the data, and the scheduled instructions can be executed
independently. Most of the out-of-order execution techniques are derived from
dataflow scheduling. However, the execution of the program is still determined
largely by the instruction sequence as dataflow scheduling is only applied to
a few instructions in a small window over the sequential code. In contrast,
Moonsoon [29] and Wavescalar [40] are based on dataflow models.
The _von Neumann_ model is simple but sequential, the _dataflow_ model is
parallel and improves the efficiency of the execution of the program but with
the cost of adding complexity in the hardware design, therefore a _hybrid
dataflow_ model [6, 17] was tried to combine the advantages of both. The issue
with generalized dataflow model is the requirement of a matching store to
detect which operations become ready to execute when its dependencies are
satisfied. The efficient organization of the matching store is not so clear to
the designers of the execution models. With a hybrid dataflow model, the
ordering of instructions can be sequentialized to reduce the need for a big
matching store. In hybrid dataflow models instructions are executed based on
the _von Neumann_ model within a thread but there is support for dynamic
execution of threads based on the _dataflow_ model. This model allowed
multicore architectures to exploit parallelism in programs. P-RISC [33],
Multiscalar [38] and DDM-CMP [39] are based on the _hybrid dataflow_ model.
Modern many-core systems one way or other provide concurrency constructs to
exploit parallelism. To quote Prof. Chris Jesshope _By 2020 we could see up to
$10^{4}$ cores and $10^{6}$ hardware threads on a single chip_. The
concurrency constructs introduced by the microthreading model are implemented
in the instruction set of the microthreaded architecture which is also
referred as the Microgrid. This architecture assumes a lot of concurrency in
the applications and provides the concurrent architecture to exploit
parallelism.
The rest of the paper is organized as follows. In Section 2 we define the
latency tolerance feature in computer architecture. We explain the details of
the microthreading model in Section 3 and the details of the architecture
based on this model in Section 4. We describe the programming of the Microgrid
in Section 5. We explain the microthreading model in the context of the
Microgrid in Section 6. We present the I/O in the Microgrid in Section 7 and
conclude the paper in Section 8.
## 2 Latency tolerance
In any program, a computation is preceded and followed by memory operation
which takes a variable amount of time as it depends on the locality of the
data i.e. the data is located in L1-, L2-,L3- caches or off-chip memory. In
single-threaded programs the processor has to wait for memory operations to
complete and then continue with the computation. In multi-threaded programs,
when a memory operation is issued, the thread may be suspended and execution
is switched to another available thread. Because of dataflow scheduling, the
memory operation completes asynchronously and wakes up the suspended thread.
The execution of program is tolerant to the long latency operations and is
shown in fig. 1 assuming multiple instructions issued and executed in a single
core.
The use of multiple threads per core in order to hide latency has been
understood for long time [53]. For instance, software time sharing is adequate
to tolerate external I/O latencies but can not tolerate other long latency
operations (e.g. memory access) because of the absence of fine-grained latency
tolerance. Hardware multi-threading with dynamic scheduling is used to
tolerate long latency operations e.g. MTA cores can switch on every miss
operation but it requires a long pipeline flush. Niagara cores switch on the
issue of a memory operation [37] and therefore tolerate long latency
operations, but can not tolerate the latency in FPUs or other asynchronous
operations (e.g. management of threads).
Figure 1: Latency tolerance in multi-threaded programs assuming multiple
instructions issued and executed in the core.
## 3 The microthreading model
The microthreading model [18] is based on a type of _hybrid dataflow_ model
and has evolved from DRISC [4] (Dynamically-scheduled RISC) which was proposed
in 1996 with the goal of separating computation from communication. DRISC
provides dataflow scheduling in RISC core, which executes instructions
asynchronously and with multiple threads it can tolerate the latency of long
latency operations. The innovation in this model is that instead of multiple
threads, an ordered set of threads referred as families are used, which
provide composability in programs. The operations to create families resemble
the fork/join operations found in most parallel programming models. The
microthreading model is actually a hybrid of von Neumann model, dataflow
scheduling of threads.
The microthreading model has been refined over the decade from a single
processor model to an abstract machine model for many-cores and can capture as
much concurrency as possible using families of threads in a dynamically
evolving concurrency tree. A family is comparable to loop or function call in
traditional programming. Any thread can create further families (any
heterogeneous combination is supported by the model) showing a hierarchy of
threads. An example concurrency tree of the microthreading model is shown in
fig. 2, where every family is composed of some number of threads and every
thread can create another family.
Figure 2: An example concurrency tree created by the microthreading model. A
family can create any combination of threads and the hierarchy of families can
go until any level.
As the overhead of creating and synchronizing threads in software is 10 to 100
thousands cycles on contemporary hardware [49]. The microthreading model
shifts the perspective from software threads to hardware threads with
concurrency management in hardware in order to reduce the cost from few
thousands cycles to just a few cycles. The number of cycles taken by an
instruction depends on the dynamic state of the architecture, but because of
asynchronous completion and fine-grained latency tolerance the model has the
potential to achieve the goal of RISC i.e. _one cycle per instruction_. In the
microthreading model the throughput of the program can potentially demonstrate
that every instruction takes one cycle to complete, assuming single
instruction issued and executed in the core.
### 3.1 Communication and synchronization
The microthreading model supports hybrid dataflow scheduling by using
I-structures [3]. An I-structure is a data structure with the semantics of
dataflow i.e. every element has a state of either _full_ or _empty_. Any
operation accessing an element with _empty_ state is suspended. The state of
the element is changed to _full_ asynchronously, and the suspended operation
is released to access the element. This process is termed as _split-phase_
operation [54] which has two phases; requesting and consuming. A request to
the desired data is issued, but if the data is not available, the request is
suspended. The execution model can continue executing other computations while
the request is in the progress. When the data arrives, the suspended operation
is released and the instruction can consume the data.
The microthreading model uses I-structure as a set of channels; globals and
shared to support communication between threads. These channels are actually
registers address using the register names of the underlying ISA. But for the
sake of generality these registers are referred as channels. These channels
have blocking read and non-blocking writes. By blocking we mean that the
operation is suspended because of the unavailability of the data, and by non-
blocking we mean that the operation is performed and the writing of the data
complete asynchronously. Global channels are mapped to all threads in the
family and have read only access. Shared channels provide the mechanism for
the communication only between adjacent threads uni-directionally in the
ordered sequence of threads in the family. The sharing of synchronizers
between threads enables fast thread-to-thread synchronization, for example to
implement dependencies in a loop. The creating thread communicate with the
created threads using messages, which can access the synchronizing memory of a
family remotely. The one-way communication is restricted but ensures a
deadlock free communication [52]. The communication and synchronization
between threads are shown in fig. 3 and more details can be found in [21, Sec.
4.3.3.3].
Figure 3: The communication between threads and the synchronization of the
family.
### 3.2 Memory consistency model
The microthreading model addresses consistency using two models for threads
and families. A single thread is sequentially consistent, as all the
instructions within a thread appear to execute in sequential order. A family
is weakly consistent, as it follows the three properties of Weak Consistency:
* •
Any created thread can perform a read or write operation only after all the
writes by the creating thread prior to the creation of the family.
* •
The family can be synchronized only when all previous write operations by the
created threads are performed.
* •
The creating thread can perform a read or write operation only after the
created family is synchronized.
Memory consistency models [11, 28] have been defined for different programming
and machine models. We present only two consistency models which are relevant
to the microthreading model.
* •
Sequential Consistency (SC) :
A system is sequentially consistent if the result of the execution of
instructions is the same as if instruction of all the cores were executed in
some sequential order. In addition, the instructions of each individual core
appear in the same sequence as specified by the program. SC has two cases:
1. 1.
A read operation by any processor/thread is allowed to perform only if all
previous read and write accesses are performed globally.
2. 2.
A write operation by any processor/thread is allowed to perform only if all
previous read and write accesses are performed globally.
* •
Weak Consistency (WC) :
A system is weakly consistent if the synchronizing operation is performed only
if there is no pending ordinary (non-synchronizing) read or write operations,
and any ordinary read or write operations can not be performed if there is a
pending synchronizing operation. The weak consistency is defined in [11] but
we are slightly modifying the terminology to match with the microthreading
model. WC has three cases:
1. 1.
All previous synchronizing operations must be performed before an ordinary
operation is performed.
2. 2.
All previous ordinary operations must be performed before the synchronizing
variable.
3. 3.
Synchronizing operations are sequentially consistent with respect to one
another.
### 3.3 The notion of resources in the microthreading model
The microthreading model is an abstract machine and programming model and
therefore does not have physical resources, but it is designed with the notion
of resources as the model addresses many-core architectures. A family of
threads is allocated to a group of resources and are referred as _place_.
These resources can be one or more cores, single or many threads etc. A family
created on a _place_ will execute there until synchronized and a thread
allocated to a core will execute there until terminated.
In order to control the number of threads on a core per family, _windowsize_
is introduced as a run-time parameter. A programmer can carefully use this
parameter to use only the requested number of resources, leaving resources for
other families (or may be more important families). More importantly this
parameter is used to avoid deadlock from dependencies down the concurrency
tree. It can wisely be used by the programmer or compiler to ensure that at
least one leaf of the concurrency tree can continue execution by using the
available resources and hence avoid deadlocks. The _windowsize_ limits the
number of threads executing in a family, similar to the concept as k-bounded
loops [8] in dataflow scheduling. The parameter ’k’ is comparable to the
_windowsize_ in the microthreading model. However, it should be noted that
_windowsize_ is per core and ’k’ is per loop which may be distributed on many
cores.
The microthreading model introduces _break_ used by a thread to terminate the
execution of the family. The _break_ will let the already created threads of
the family to complete and will stop the creation of any new threads in the
family. It is comparable to the _exit_ statement in loops or functions in
traditional programming, and is useful when the synchronization of the family
depends on some dynamic conditions.
### 3.4 Concurrency constructs
The microthreading model defines concurrency constructs to use the concurrency
provided by the hardware. But these constructs need to be introduced in some
intermediate language in order to write programs. In this section we explain
some of the concurrency constructs of the microthreading model and show that
these constructs are defined in an intermediate language referred as SL [21,
22]. In Section 5 we will show an example program that will demonstrate the
use of concurrency constructs of the microthreading model in an application.
### Allocation and deallocation of place
A group of cores i.e. _place_ is required to be allocated before the creation
of a family and released when the family is synchronized.
* •
$sep\\_alloc(<parameters>)$ and $sep\\_free(<parameters>)$
are used to allocate and de-allocate cores through a software API. The
$<parameters>$ are a list of parameters passed for allocation and de-
allocation e.g. specifying the strategy of allocation or number of cores etc.
and can be found in [19].
### Creation and synchronization of family
A family of threads can be created by using $sl\\_create$ constructs similar
to a loop in traditional programming. The created family can be synchronized
using $sl\\_sync$ construct.
* •
$sl\\_create(fid,pid,start,limit,step,windowsize,options,thread[,<arguments>])$
defines the creation of a parameterized family.
* •
$fid$
is the identifier of the family of type $sl\\_family\\_t$.
* •
$pid$
is the identifier of the _place_ of type $sl\\_place\\_t$.
* •
$start$, $limit$ and $step$
indicate the starting, ending and iteration step counter of threads (similar
to the iterators in a loop).
* •
$windowsize$
determines the upper limit on the number of threads that can be created on a
core.
* •
$option$
decides the way the family should be created an executed. e.g. $sl\\_forceseq$
will force the family to execute sequentially.
* •
$thread$
is the name of the thread which defines the code to be executed by created
threads in the family.
* •
$<argument>$
is a comma separated list of global and shared arguments passed to the thread.
* •
$sl\\_sync()$
is used for the synchronization of the family.
### Global and shared channels
The global and shared channels introduced by the microthreading model are
implemented as global and shared variable in the intermediate language. They
are used as parameters/arguments to the creation of a family of thread.
* •
$sl\\_glarg(type,variable,value)$ and $sl\\_sharg(type,variable,value)$
indicates the global and shared $variable$ of $type$ with some $value$ passed
from the creating thread to the created thread.
* •
$sl\\_glparm(type,variable)$ and $sl\\_shparm(type,variable)$
defines the global and shared $variable$ of $type$ received by the created
thread from the creating thread.
* •
$sl\\_getp(variable)$
is used to read from global or shared $variable$.
* •
$sl\\_setp(variable,value)$
is used to write the shared $variable$ with a modified $value$ to be read by
the next thread.
* •
$sl\\_seta(variable,value)$
is used by the creating thread to write a $value$ to the shared variable to
the first thread of the family.
* •
$sl\\_geta(variable)$
is used by the creating thread to read the shared $variable$ from the last
thread in the family.
### Starting and terminating threads
Threads are implicitly created with the creation of a parameterized family and
terminated when all of their instructions are executed. But programmers need
to define the starting and terminating of a thread similar to the way
functions are defined in traditional programming.
* •
$sl\\_def(type,<parameter>)$
is used to define a thread with return of $type$ and $<parameter>$ of a list
of the global or shared variables used by the thread.
* •
$sl\\_enddef$
defines the terminating of the thread.
* •
$sl\\_index(index)$
can be used to retrieve the $index$ of the ordered set of threads in the
family.
### Breaking family
Any thread in a family can use the break statement to stop creation of new
threads in the family. The syntax is given below:.
* •
$sl\\_break$
is used to stop creating new threads in the family.
### 3.5 More about families
The microthreading model has evolved over the years and have introduced
different types of families as per the requirements of the programs. All the
families are created using $sl\\_create$ construct (except detached family,
see below). The family can be of any of the type (described in this section)
based on the way the created threads are executing. However, some types of
families are need to be explicitly defined by passing a parameter to the
create construct. In this section we give a brief overview of the different
types of families supported by the microthreading model with the aim that
programmers can write a parallel program using different combination of
threads and families. Families that are required to be explicitly defined at
the time of creation are explicitly stated in their respective subsection.
#### Independent family
A family is called independent family when its created threads do not require
any communication between threads. Embarrassingly parallel applications
generally do not require communication between created threads and the
microthreading model supports these types of applications by creating them as
independent families.
#### Dependent family
Threads in dependent family communicate with each other. Fine-grained parallel
applications require a lot of communication and coarse-grained parallel
applications communicate a little between threads. The microthreading model
supports these types of application by creating dependent families. These
families communicate with each other through shared channels introduced
earlier (c.f. Section 3.1). Dependent families are inherently sequential and
therefore do not get any speedup by distributing the threads on many cores.
But in the microthreading model these threads get benefit from latency
tolerance and asynchronous completion as a dependent family can be executed
along with other families.
#### Homogeneous family
All families are statically homogeneous but can be made dynamically
heterogeneous if the programmer calls/creates different functions in a thread
dependent on the index of the thread.
#### Heterogeneous family
A family can be composed of threads that are not identical to each other. The
programmer need to explicitly write the code for different threads based on
the index of the thread to perform different operations.
#### Detached family
A family is required to be synchronized before continuing further. But there
may be situations when a family may not want to synchronize and is called
_detached_ family. For instance, creating a family that prints some characters
on terminal may not require synchronization as all the created threads perform
only read operations and do not modify any memory. The programmer has to
explicitly specify the _detached_ family by creating a family with
$sl\\_detach(<parameters>)$ instead of $sl\\_create(<parameters>)$, where
$parameters$ are defined in Section 3.
#### Exclusive family
To support mutual exclusion in the microthreading model, exclusive families
are introduced. Exclusive families are treated differently to regular families
because of the requirement of being mutually exclusive. More details about
mutual exclusion in the microthreading model can be found in [21, Sec. 14.1].
To create an exclusive family the programmer have to explicitly use
$sl\\_exclusive$ in the place of $options$ in the create construct. Generally
exclusive families consist of a single thread.
#### Sequentialized family
A family can be forced to execute sequentially, to avoid using resources by
creating many threads. The programmer has to explicitly use $sl\\_forceseq$ at
the place of $options$ in the create construct. The sequentialized family will
execute as regular code of the parent thread.
## 4 The Microgrid
Microgrid is actually a blueprint with parameters defined for the
architecture. There are many Microgrids; with 1 core, 128 cores, 1000 cores,
random bank memory, COMA etc. We use the term Microgrid to generalize all
instances of the blueprint. At some time we talk about a particular instance,
but then we explicitly give the details of the parameters. In the context of
the Microgrid, a core is explicitly termed as microthreaded core to
differentiate from traditional cores. In addition, a thread is termed as
microthread to distinguish it from traditional threads.
The basic component of the Microgrid is the microthreaded core, and it is
important that we understand the execution of instructions in the core in
order to simulate the core at the high-level. The concurrency constructs of
the microthreading model are implemented in the instruction set (ISA) of the
core. The cores are designed to be simple; single issue, in-order, 6-stage
pipeline [20] and based on RISC. Some components of the operating system are
implemented in the silicon of the core [7] e.g. scheduler, resource
allocation, mapping etc. Some energy inefficient features are removed from the
design of the core e.g. speculation, out-of-order execution and cache
prefetching.
Figure 4: The microthreaded core.
The 6 stage pipeline of the core is shown in fig. 4, where we are not showing
the scale of components, but mainly the layout of components. It is a classic
RISC core with some modifications to existing components and adding some more
components to support the concurrency constructs of the microthreading model
in the core. The instruction fetch (IF), instruction decode (ID), register
read (RR), execute (EX), memory (MEM), write back (WB), I-cache, D-cache,
Integer register file, Floating register file and shared FPU are existing
components but slightly modified in order to support the ISA of the Microgrid.
The scheduler, allocator, network, thread table and family table are the newly
added components to support the concurrency constructs of the microthreading
model. The thread table and family table are used to store the thread contexts
and family contexts. Every core also has a single exclusive context for
creating an exclusive family and thread. The shared channels in the
microthreading model are implemented in the registers of the cores. These
registers are synchronizing and provide two ports for synchronous and
asynchronous completion e.g. D-cache and FPU operations are completed
asynchronously and therefore connected to the asynchronous ports of the
registers. The bus interface connects the L1-cache to the snoopy bus and the
network interface connects the core to the network of other cores on the chip.
The threads in hardware are supported by the thread management unit through
the instruction set of the core for creation and synchronization of threads.
It has lowered the overhead of creation and synchronization of software
threads from more than ten thousands of cycles to just few cycles. It also has
extremely low overhead of context switching (zero cycle) and provides fine-
grained interleaving i.e. interleaving at every cycle. However, this
interleaving is not as strict as in traditional threads. Threads can execute
multiple instructions from a single thread until a context switch is required.
Interleaving can be bypassed to maximize pipeline usage if there is no other
available thread. In single-threaded programs interleaving does not make sense
anyway. To ensure fairness in the execution of threads, interleaving is
enforced so that no thread can monopolize the execution time of the core.
The binary code of a program generated for the microthreaded core can be
executed on any other core that supports the same instruction set as
Microgrid. The other core may not support concurrent execution and the code
will be executed as single threaded. The code can also be executed on
different number of thread slots per core. This is called binary code
compatibility i.e. the binary code may not get the same performance but at
least it will not fail, its performance will be gracefully degraded.
As Microgrid provides many cores on a single chip, we show a group of four
cores in fig. 5 to demonstrate the interconnection of large number of cores.
The delegation network is a Network-on-Chip (NoC) where all cores are
addressable from all other cores. Currently it is implemented as a fully
connected network in MGSim, but could be implemented as a mesh network in the
hardware implementation. A core is connected to the previous and the next core
by a distribution network. Two cores share an FPU, and every core has an
L1-cache which are connected to L2-cache by a snoopy bus. L2-caches are
connected with each other in the distributed cache network. In order to handle
deadlock in memory network, every L2-cache has an incoming and outgoing
message buffer. An example layout of the 128-core Microgrid in a single chip
is shown in fig. 6.
Figure 5: A group of four cores, each core has an L1-cache, two cores share an
FPU and four cores share an L2-cache. Figure 6: The layout of 128 cores on the
Microgrid chip.
### 4.1 The communication network
The Microgrid chip has two communication networks on the chip, in addition to
the distributed memory network.
#### Delegation network
Every core is connected to every other core by a bi-directional fully
connected network. This network is highly efficient because a message travels
in 10 or even less cycles from source to destination. The current
implementation of the Microgrid assumes a single cycle, node-to-node routing
in a lightly loaded network, but in a loaded network it may take more cycles
due to contention and buffering delays. Since delegation network is very
efficient, it is used only for the concurrency management between cores that
are not adjacent.
#### Distribution network
All the cores are connected in a single bi-directional daisy-chained network
in a Moore curve. The curve can be chosen in a way to preserve locality in
L1-caches and L2-caches. It takes two cycles to travel from one core to an
adjacent core. The cycles taken by a return trip of a message from a core to
any other core on the distribution network is given in Eq: 1, where $c$ is the
number of cores.
$delay=2\times 2\times c$ (1)
The distribution network is used for the logical partitioning of the chip i.e.
different parts of the program can be executed on different parts of the chip.
It provides grouping of cores to be used for delegating a family, where the
distribution network knows the starting core in the group and the size of the
group.
### 4.2 Resource management in the hardware of the Microgrid
An integer value in the source program is used to identify a group of adjacent
cores where a family can be delegated. We refer to the group of cores as
_place_ and the integer value that identify the place as _placeid_. The
identifier can identify the starting core of the place and the size of the
place. More details about places can be found in [24, 22] and [21, Chap. 11 &
App. E].
A family can execute on the same core as the parent thread or the same place
as the parent thread. The former is called _local_ place and the later is
called _default_ place and are identified by place id as $0$ and $1$
respectively. The execution of a family on a place different than used by the
parent thread is called _delegation_. We can perform some arithmetic
operations on _placeid_ to derive the starting core and the size in the place.
The starting core can be calculated as given in Eq: 2 and the size of the
place can be calculated as given in Eq: 3. The size of the place is always in
power of 2, and the starting core must have the _coreid_ which is multiple of
the size of the place. This is a restriction, but easy in implementation as
this information can be derived from integers.
$starting\\_core\;\;=\;\;(placeid\;\;\&\;\;(placeid\;-\;1))\;>>\;1$ (2)
$size\;\;=\;\;placeid\;\;\&\;\;-placeid$ (3)
### 4.3 Proposed software service to access resources in the Microgrid
A software layer for the allocation of resources is defined in a protocol
referred as SEP [19, 27] which provides an easy to implement and efficient
management of resources. In the initial research work, a single core on the
Microgrid is reserved for the operations of SEP. It works in mutually
exclusive manner, and can soon become a bottleneck when a lot of requests are
coming to this single core. Some research is ongoing in avoiding bottleneck
and making the management hierarchical.
SEP is implemented using the well-known binary buddy allocation [30] of the
memory management. An example of buddy allocation used for cores in the
Microgrid is shown in fig. 7. It shows a system having a total of eight cores
with their identifiers given on the top. The time the system is initialized,
all the cores are in one group. Then a request is issued to allocate one core
and involves these steps:
* •
Divide 8 cores in half.
* •
Still bigger than required, divide 4 cores into half.
* •
Still bigger than required, divide 2 cores into half.
* •
We found a single core as requested, allocate the single core.
A request for a given number of cores, will allocate the group of cores if
available. In case the available group is larger than requested, divide the
bigger group into smaller group until the group of requested cores is
obtained. When the allocated cores are released, they are grouped into a
larger group. In the given example when the last 4 cores are released it is
grouped into larger group of 8 cores, bringing the system to the initial
state.
Figure 7: An example of the binary buddy allocation of cores by SEP.
A programmer can choose different numbers of cores in a place. Every time a
place is requested, a group of cores which are not allocated to any other
family is allocated. The allocated cores can be used by the family (and sub
families if any) to which the place is allocated, and can not be allocated to
any other family until explicitly de-allocated. SEP provides two API to
programmers; $sep\\_alloc()$ and $sep\\_free()$ for allocation and de-
allocation respectively similar to $malloc$ and $free$ in C. The allocation
and de-allocation of SEP can be performed asynchronously, but in the current
implementation these processes are not asynchronous. The moment the allocation
or de-allocation is issued by a thread, the thread will wait until the
operation is completed. The allocation of place depends on the policy given
below:
* •
Minimum: Allocate at least the number of cores specified (may be more if
available).
* •
Maximum: Allocate at most the number of cores specified (may be less if
unavailable).
* •
Exact: Allocate the exact number of cores. The allocation will fail if the
given number of cores can not be allocated.
* •
Any size: Any available size in the power of 2 will be allocated, starting
from the lowest available size of the group.
## 5 Programming the Microgrid
In this section we show an example program i.e. Matrix Multiplication of equal
sized matrices, to demonstrate the way programs are written for the Microgrid.
The objective is to show that a sequential C program can be transformed to a
microthreaded program easily by high-level programming languages, a compiler
or with little effort by the programmer. We also show the concurrency
constructs in the generated assembly to demonstrate the way the concurrency
constructs are supported in the ISA of the Microgrid.
#### Sequential C program
The sequential C code for Matrix Multiplication of size $1000\times 1000$ is
shown in listing LABEL:code:seq_matmul. We allocate three arrays; two for
source matrices and one for the result matrix. Once the memory is allocated,
we can fill these arrays with some numbers, but for saving space in the page,
we assume the existing values present in those memory locations. To perform
matrix multiplication we write three loops; outer loop, middle loop and inner
loop. The inner loop is the one which multiplies the two elements of the
source matrices and stores in the result matrix. After the multiplication is
completed, we free the allocated memory.
Listing 1: Sequential execution of Matrix Multiplication.
⬇
#include <stdio.h>
int N = 1000;
int main()
{
// Allocate memory for matrices
int *A = (int*) malloc(N * N * sizeof(int));
int *B = (int*) malloc(N * N * sizeof(int));
int *C = (int*) malloc(N * N * sizeof(int));
int i, j, k;
// Perform the multiplication
// matmul_outer
for(i=0; i<N; i++)
{
// matmul_middle
for(j=0; j<N; j++)
{
//matmul_inner
for(k=0; k<N; k++)
{
C[i*N+j] += A[i*N+k] * B[k*N+j];
}
}
}
free(A), free(B), free(C);
return 0;
}
#### Microthreaded program
The microthreaded program for the Matrix Multiplication of size $1000\times
1000$ written in SL is shown in listing LABEL:code:par_matmul. In the
$t\\_main$ function, we first allocate a group of 8 cores in the Microgrid.
Then we allocate three arrays for source and result matrices. Then we replace
the outer loop by creating a family of $N$ threads. The threads in the outer
family create further middle families and then each of those creates the inner
family. The inner family perform the multiplication of elements in source
arrays and store the multiplication in the resultant array. The transformation
from sequential program to microthreaded program involves the creation of
families and shared and global parameters. Once all the families are
synchronized, the allocated memory to arrays and allocated cores are released.
Listing 2: The microthreaded version of Matrix Multiplication.
⬇
#include <svp/sep.h>
#include <stdio.h>
int N = 1000;
//matmul_inner
sl_def(matmul_inner, void,
sl_shparm(long, sum), sl_glparm(void*, A),
sl_glparm(void*, B), sl_glparm(size_t, i),
sl_glparm(size_t, j)) {
sl_index(k);
int (*A)[N][N] = (int (*)[N][N])(void*) sl_getp(A);
int (*B)[N][N] = (int (*)[N][N])(void*) sl_getp(B);
int v = (*A)[sl_getp(i)][k] * (*B)[k][sl_getp(j)];
sl_setp(sum, v + sl_getp(sum));
}sl_enddef
//matmul_middle
sl_def(matmul_middle, void,
sl_glparm(void*, A), sl_glparm(void*, B),
sl_glparm(void*, C), sl_glparm(size_t, i)) {
sl_index(j);
sl_create(,,0,N,,,,matmul_inner,
sl_sharg(long, sum, 0), sl_glarg(void*, , sl_getp(A)),
sl_glarg(void*, , sl_getp(B)), sl_glarg(size_t, , sl_getp(i)),
sl_glarg(size_t, , j));
sl_sync();
int (*C)[N][N] = (int (*)[N][N])(void*) sl_getp(C);
(*C)[sl_getp(i)][j] = sl_geta(sum);
}sl_enddef
//matmul_outer
sl_def(matmul_outer, void,
sl_glparm(void*, A), sl_glparm(void*, B),
sl_glparm(void*, C)) {
sl_index(i);
sl_create(,,0,N,,,,matmul_middle,
sl_glarg(void*, , sl_getp(A)), sl_glarg(void*, , sl_getp(B)),
sl_glarg(void*, , sl_getp(C)), sl_glarg(size_t, , i));
sl_sync();
}sl_enddef
sl_def(t_main, void) {
// Allocate a place using SEP
int core = 8;
sl_place_t pid;
if (sep_alloc(root_sep, &pid, SAL_EXACT, core) == -1) {
printf(”cannot␣allocate␣a␣place.”);
exit(1);
}
// Allocate memory for matrices
int *A = (int*) malloc(N * N * sizeof(int));
int *B = (int*) malloc(N * N * sizeof(int));
int *C = (int*) malloc(N * N * sizeof(int));
// Perform the multiplication
sl_create(,pid,0,N,,,,matmul_outer,
sl_glarg(void*, , A), sl_glarg(void*, , B), sl_glarg(void*, , C) );
sl_sync();
free(A); free(B); free(C); sep_free(root_sep, &pid);
} sl_enddef
#### Microthreaded assembly code
The microthreaded assembly code is shown in listing
LABEL:code:assembly_matmul. We can see the assembly instructions of
concurrency constructs i.e. $allocate$ for the allocation of family, $crei$
for creation of family, $sync$ for synchronization of family, $puts$ writing a
shared variable, $gets$ reading from a shared variable etc. This is the
innovation in the Microgrid, which avoids the mapping of software threads to
hardware threads, which reduces the overhead from more than 10-100 thousands
cycles of software threads to just few cycles in hardware threads. The ISA of
program is scheduled by the microthreaded core with the hardware support for
concurrency management.
Listing 3: A snipper of the assembly program generated from the microthreaded
program of Matrix Multiplication for the Microgrid.
⬇
…
# Assembly code of thread matmul_inner
.ent matmul_inner
.registers 4 1 4 0 0 0
matmul_inner:
ldpc $l3
ldah $l3, 0($l3)
lda $l3, 0($l3)
$matmul_inner..ng:
ldah $l1,N($l3)
ldl $l1,N($l1)
mulq $l0,$l1,$l2
mulq $g2,$l1,$l1
addq $l2,$g3,$l2
addq $l1,$l0,$l0
s4addq $l2,$g1,$l2
s4addq $l0,$g0,$l0
ldl $l2,0($l2)
ldl $l0,0($l0)
mull $l2,$l0,$l0
addq $l0,$d0,$l0
mov $l0, $s0
end
.end matmul_inner
…
# to show the concurrency constructs in the ISA
.ent t_main
…
allocate/s $l0, 0, $l0
setstart $l0, 0
ldq $l1,matmul_outer($l17)
setlimit $l0, $l9
setstep $l0, 1
setblock $l0, 0
wmb
crei $l0, 0($l1)
putg $l12, $l0, 0
putg $l11, $l0, 1
putg $l10, $l0, 2
sync $l0, $l1
mov $l1, $31
release $l0
…
.end t_main
#### SL tool-chain
Figure 8: The SL tool-chain.
The SL tool-chain is shown in fig. 8 and is built around SL [22] which is a C
based language but is extended to express the concurrency constructs of the
microthreading model. SL is an intermediate language intended for higher level
programming languages (such as Single Assignment C or SAC [36, 13, 12, 14,
15], FastFlow [2, 1] etc.) and parallelizing C compiler [34, 35]. SAC is
Matlab-like programming language, and it provides concurrency from a very high
view. Programmer writes program using arrays and let the compiler decide to
exploit concurrency in the programs.
A number of tools and simulators are added to the designer’s toolbox and used
for the evaluation of the Microgrid from different perspective. The SL
compiler can generate binary for different implementations of the Microgrid.
We have software libraries that provide the run-time systems for the
microthreading model on the shared memory SMP machines and referred as _svp-
ptl_ [51] and distributed memory for clusters/grids and are referred as Hydra
[26] and _dsvp-ptl_ [50] The SL compiler can generate binary for UTLEON3 [9,
10], MGSim [5, 32] and HLSim [44, 45, 46, 43, 47, 48, 41, 42].
Unless specified otherwise, the SL compiler generates two implementations for
every family in the given SL program; concurrent and sequential. Either one is
used as per the dynamic state of the chip. In case resources can not be
allocated to a family, the sequential version can be used. The sequential
version of the program is used to avoid deadlock and if the deadlock can be
avoided by static analysis, then suspension on resource allocation is safe.
The programmer can force the program to use the concurrent version all the
time i.e. threads are suspended until resources become available.
## 6 The microthreading model in the context of the Microgrid
We gave details of the Microgrid, and would like to revisit the microthreading
model in the context of the Microgrid. We need to understand all the details
of the model and the architecture in order to simulate the architecture at a
high level, as the high-level simulator have to exhibit the same behavior.
### 6.1 A concrete example
Figure 9: A concrete example of the microthreaded program of the
microthreading model executing in the single core of the Microgrid. We
explicitly use a single core to show the coordination of different events in
different threads over time.
A concrete example of the microthreaded program in the microthreading model
executing in a single core of the Microgrid is shown in fig. 9. All
concurrency constructs are shown in the figure in the form of rectangles of
different sizes/color. We also show the effect of concurrency constructs on
the execution of threads. Every thread has an implicit _start_ and _end_
event. Threads are created at the rate of one thread per cycle. The creation
process takes 4 cycles (with additional latency to load the cache-line from
memory if not in the I-cache, also potential queueing delays if another family
is being created). The instructions of a thread can be executed as soon as the
thread is created which means the first instruction of the first thread can
execute along with the creation of the other threads [23]. The process of
creation of threads is implemented sequential on a single core, and is decided
based on the trade-off between simple design and efficient creation of family.
In the figure we can see that two _allocate_ events execute concurrently, but
the _create_ events are sequentialized on a single core.
We show a heterogeneous family of threads i.e. Threads $0$ and $1$ in Family
$0$ in the first level is different from each other. We show a dependent
family where threads are waiting for the previous threads to write the shared
variable/register. We also show that the sync event is waiting for all the
threads to complete. We show a homogeneous independent family where the parent
threads is executing along with the child threads and therefore the sync does
not wait as the time the sync is issued where all the created threads are
completed already.
### 6.2 Family’s Life cycle
A family passes through different stages in the Microgrid during the
execution. In this section we consider an example when a parent thread is
executing on a core and it delegates a family to a place of four cores shown
in fig. 10. A detailed explanation of these stages is given in below section.
Figure 10: The life cycle of a family during the execution in the Microgrid.
#### Allocation
It is implemented as try-to-allocate mechanism (see below for various options
on allocate). In this process; family context, thread contexts and registers
on all cores of the place are tried to be allocated. The parent thread sends a
message on the delegation network to the first core. The first core checks the
availability of at least one thread context, one family context and 31
registers. In case of success the message is passed to the next core of the
place using distribution network. When all the cores succeed, the message
comes back from the last to the first core on distribution network, and the
contexts are asynchronously allocated. The first core then notifies the parent
core using delegation network. The allocate protocol has three different
modes:
* •
Normal (or default): The allocation fails immediately after the allocation
fails on any core in the place.
* •
Suspend: The allocation waits for the availability of resources.
* •
Exclusive: The allocation allocates exclusive contexts and in case of
unavailability, keeps waiting until resources are acquired.
The allocate protocol has four different strategies:
* •
Normal (or default): It will try to allocate as many cores as requested,
possibly down to one core (in power of 2).
* •
Exact: It will allocate the family on all the requested cores.
* •
Single: It will allocate the entire family to the first core in the place.
* •
Balanced: It will allocate the entire family to the least loaded core of the
place, in terms of family contexts.
#### Configuration
Once the contexts are allocated on the cores, the next message is sent on the
delegation network to the first core to store parameters in the allocated
family table entry for the bulk creation. It consists of $setstart$,
$setlimit$, $setstop$, $setblock$ etc. which are derived from the create
construct and defines the number of threads to be created per core. The
configuration process is completed asynchronously and does not need any
acknowledgement. Therefore as soon as the configuration message is issued, the
next message can also be sent. In the current implementation, the delegation
network preserves the order of messages i.e. configure messages will reach to
the destination before create message.
#### Creation
The create message is sent from the parent core to the first core in the place
via delegation network. As soon as the message is received by a core, thread
creation is started asynchronously and the message is forwarded to the next
core in the place. The first core will acknowledge the parent core when all
the threads on the first core are created. The parent core may then write the
global and/or shared registers to the first core through inter-context
communication and can continue with its instructions (if any) or wait for the
synchronization of the family.
When all the threads are created, the core can de-allocate the contexts that
are not really required e.g. at the time of allocation 31 registers were
allocated, but suppose a thread is created which only uses 15 registers
therefore the core will de-allocate the extra 16 registers. All the threads
continue their execution on the allocated core, and as the threads get
terminated the context get de-allocated asynchronously. Every core shares the
information when it completes executing the allocated threads of a given
family with the next core. The last core has the information that all the
threads in the family are completed.
#### Synchronization
When all threads in the family complete, the last core will send a message on
the delegation network to the parent core. The parent thread is activated so
that the family can be synchronized. All the modified state from the family
becomes defined. In the case of a dependent family, the parent core will also
read the value from the register updated by the last thread through inter-
context communication.
#### Release
A message is sent from the parent core to the first core on delegation network
to release the family. This process will complete asynchronously and the
parent core can continue its execution.
### 6.3 Thread’s life cycle
Figure 11: The lifecycle of a thread during the execution on the core of the
Microgrid.
A thread passes through different stages during the execution and is shown in
fig. 11. A newly created thread is allocated thread context, registers and PC
(Program counter) is configured to place the thread in ready queue. In case
the thread gets an I-cache miss the thread waits until the instructions are
loaded from cache. The thread then passes to the active queue. The context of
the thread is moved from the active queue to running where the instructions
actually execute in the execute stage of the pipeline. A thread can be
suspended when it is active or running at the write back stage of the pipeline
e.g. in case of dependency. The asynchronous completion of a thread moves the
thread from the suspended stage to ready queue. At the termination of the
thread the allocated entries are cleaned up which can be used by any newly
created threads.
### 6.4 Distribution of threads
Figure 12: Distribution of independent threads on 4 cores with 5 window size.
The distribution of threads to the cores in a place depends on the type of
family i.e. dependent or independent. The threads in a dependent family are
restricted to single core only, as no significant speedup can be obtained by
distributing them on multiple cores. However instructions in the thread will
get benefit from latency tolerance on asynchronous completion in loads, store
and floating point operations.
Threads of an independent family are distributed by an equal distribution of
threads per core. In fig. 12 we show an example of a family with 40 threads on
4 cores with window size of 5. Threads are distributed as 40/4 and every core
can create 10 threads. But the window size is 5 and therefore only 5 threads
can be activated at a time on a given core. As soon as a thread terminates, a
new thread can use the context of the terminated thread. For instance on core
x, threads 0,1 and 2 are terminated, threads 3-7 are executing and threads 8
and 9 are still waiting to be created.
### 6.5 Communication through registers
Figure 13: Allocation of registers to threads and their communication through
registers.
Every thread created in the Microgrid uses a set of registers in the register
file. The channels of the microthreading model used for the communication and
synchronization of the family introduced in Section 3.1 are implemented in
registers of the Microgrid. The registers allocated to a thread are
categorized as; globals, locals, shareds and dependents. The mapping of
registers to threads and their communication through registers is shown in
fig. 13. Some explanation of these registers are given as:
* •
Global registers implements global channels and are visible to all threads in
the family. The parent thread writes to the global registers and all created
threads in the family can read from them.
* •
Local registers are only visible to the individual thread only.
* •
Dependent registers implements shared channels and have read-only access to
the shared registers from the previous threads.
* •
Shared registers implements shared channels and have the write access by the
current threads, and carry the modified value to be read by the next thread.
## 7 I/O in the Microgrid
The Microgrid requires a decentralized approach where not every core is
connected to I/O but some specialized microthreaded cores are used to support
I/O and referred as I/O cores. I/O cores have limited/extended instruction set
[25] compared to the regular microthreaded core and has no floating point
operations. The instruction set of the I/O core supports the I/O
infrastructure and are connected to the delegation network of the Microgrid.
We do not simulate I/O cores in the current implementation of HLSim, therefore
we refer readers to [16, 21, 31] for the I/O management in the Microgrid.
## 8 Conclusion
The microthreading model is a hybrid dataflow model and provides the
simplicity of _von Neumann_ model and asynchronous completion of the
_dataflow_ execution model. It is an abstract machine model for many-cores
architecture and shifts the perspective from software threads to hardware
threads. Because of asynchronous completion and fine-grained latency tolerance
the microthreading model can potentially achieve the goal of RISC i.e. _one
instruction per cycle_ in the throughput of the program assuming single issue
width. The details of a particular type of future many-cores systems are given
to demonstrate the complexity of the architecture.
## Acknowledgement
The author would like to thank Dr. Raphael Poss, Dr. Michiel van Tol and Prof.
dr. Chris Jesshope.
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|
arxiv-papers
| 2013-09-21T17:50:09 |
2024-09-04T02:49:51.286855
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Irfan Uddin",
"submitter": "Irfan Uddin",
"url": "https://arxiv.org/abs/1309.5507"
}
|
1309.5551
|
# Design space exploration in the microthreaded many-core architecture
Irfan Uddin
University of Amsterdam, The Netherlands
[email protected]
###### Abstract
Design space exploration is commonly performed in embedded system, where the
architecture is a complicated piece of engineering. With the current trend of
many-core systems, design space exploration in general-purpose computers can
no longer be avoided. Microgrid is a complicated architecture, and therefor we
need to perform design space exploration. Generally, simulators are used for
the design space exploration of an architecture. Different simulators with
different levels of complexity, simulation time and accuracy are used.
Simulators with little complexity, low simulation time and reasonable accuracy
are desirable for the design space exploration of an architecture. These
simulators are referred as high-level simulators and are commonly used in the
design of embedded systems. However, the use of high-level simulation for
design space exploration in general-purpose computers is a relatively new area
of research.
###### Contents
1. 1 Introduction
2. 2 Design space exploration in embedded systems
3. 3 Application -dependent and -independent DSE
4. 4 Design space exploration in general-purpose systems
5. 5 Design space exploration in the Microgrid
6. 6 Conclusion
## 1 Introduction
Simulators with high simulation speed and less complexity are desirable for
early Design Space Exploration (DSE) of the architecture. Any decision to
improve the architecture becomes more expensive and requires more effort at
the later stage and requires more effort, time and budget. DSE is performed in
all kinds of computer systems. However, in the embedded systems domain the use
of high-level simulation for DSE purposes has been accepted as an efficient
approach for more than a decade. In that sense, the DSE in embedded systems
pioneered the high-level simulation techniques. Therefore, in this paper we
will give details about DSE in embedded systems using high-level simulators.
We will explain that general-purpose computers are getting more complex and
therefore high-level simulators are also required for the DSE. Since Microgrid
is a complex architecture therefore we need to know its design space before we
present the high-level simulation techniques for the Microgrid.
The rest of the paper is organized as follows. In Section 2 we give an
explanation of DSE in embedded systems. In Section 3 we differentiate the
application dependent DSE and independent DSE. We give the DSE for the
Microgrid in 5 and conclude the paper in 6.
## 2 Design space exploration in embedded systems
Embedded systems perform predefined tasks and therefore have particular design
requirements. They are constrained in terms of performance, power, size of the
chip, memory etc. They generally address mass products and often run on
batteries, and therefore need to be cheap to be realized in silicon and power
efficient. Modern embedded systems, typically have a heterogeneous
MultiProcessor-System-on-Chip (MP-SoC) architecture, where a component can be
fully programmable processor for general-purpose application or a fully
dedicated hardware for time critical applications. This heterogeneity makes
embedded systems more complex and therefore designers use high-level simulator
to perform DSE at an early stage, because high-level simulators take less
effort to develop and less time in executing applications. In this section we
describe the high-level simulation technique used for the DSE in embedded
systems.
Different high-level simulation techniques are introduced for the DSE in
embedded systems, and are often based on the separation of concerns [26]
between application, architecture and mapping functions. DSE in embedded
systems is generally application-dependent or scenario-based. Traditional
embedded systems are targeting one particular architecture and application,
the aim is to explore the design for improvement based on certain objective.
Scenario-based DSE [43] is the process of mapping every individual process of
an application to every architecture component with different configurations.
This mapping results in an exponential number of mapping choices i.e. design
space. We show an example in fig. 1 (taken from [21]) to demonstrate that only
three processes are mapped to three architecture components, but the resulted
design space is large.
Figure 1: The mapping of three application processes to three architecture
components resulting into a large number of design space to be explored.
Ideally DSE would like to considers all possible mappings, but an exhaustive
search is infeasible. Therefore computer architects use design pruning to
optimize the search through the design space to speed up the DSE. A smart DSE
intelligently evaluates a small fraction of the design space to come up with a
sub-optimal solution. These choices have a crucial impact on the success of
the final product. DSE addresses multiple objectives [15] e.g. maximum
performance, minimum power consumption and less complex components. It is very
difficult to have a single solution that meets all the objectives
simultaneously. The main problem is that the objectives are conflicting e.g.
low power generally means bad performance or good performance means high power
usage. Therefore a set of solutions are selected based on a Pareto optimal
front [1], where solutions are not dominated by any other solution looking for
the same objectives.
### 2.1 Related work
High-level simulators have been used for the DSE in embedded systems domain
for more than a decade, and are used in the research of academia and
industries. Below are some of the research groups using high-level simulation
for DSE in embedded systems. There might exist other areas of research in
using high-level simulation for the DSE of embedded systems.
* •
Sesame, University of Amsterdam [16].
* •
(Metro)Polis, University of California, Berkeley [50].
* •
Mescal, University of California, Berkeley [12].
* •
Milan, University of Southern California, Los Angeles [2].
* •
The octopus toolset, University of Eindhoven [4].
* •
SystemC-based environment, STMicroelectronics [47].
## 3 Application -dependent and -independent DSE
We want to clearly distinguish between application-dependent DSE in
traditional embedded systems and application-independent DSE in modern
embedded system or general-purpose computers. In traditional embedded systems,
applications are statically mapped to different configurations of an
architecture using some mapping functions. Based on the simulation results,
innovative ideas can be generated which can improve application, mapping and
architecture separately.
In modern embedded systems we do not have one particular application or
scenario, but a range of applications targeted to a different configurations
of the architecture. For instance in smart phones it is not only one type of
application that can statically be mapped, but a range of different types of
applications are required to be explored on the different configuration of the
architecture. In a way modern embedded systems are converging to the general-
purpose systems. The range of applications increases in general-purpose
computers, where a variety of applications can be executed on the given
architecture. In these situations, the mapping of the application to
architectural component can not be analyzed statically but instead the code
patterns in algorithms are analyzed, and then different processes of an
application are dynamically mapped to different parts of the chip based on
certain objectives. Because of the dynamic mapping, application-independent
DSE is not as trivial as scenario-based DSE.
Design pruning is more structured in traditional embedded systems. For
instance, genetic algorithms, simulated annealing etc. are some of the
structured techniques that are commonly used in design pruning. However, for
design pruning in modern embedded systems or general-purpose systems, there
exists no structured solution that can dynamically determine a reduction in
the design space to optimize the search.
## 4 Design space exploration in general-purpose systems
The growing number of cores and size of the on-chip memory are creating
significant challenges for evaluating the design space of future general-
purpose computers. We need scalable and fast simulators for the exploration of
large number of cores on a chip within limited development time and budget.
Commercially available processors available in the market have few cores on a
chip e.g. Intel E708800 Series, IBM’s POWER7 and AMD’s Opteron 600 Series. In
the near future we believe there will be hundreds of cores per chip and DSE at
the early stage can no longer be avoided [7] in general-purpose computers, as
the number of mapping an application explodes as the number of cores
increases.
The use of high-level simulators for the DSE in general-purpose computers is
relatively new compared to embedded systems domain. A number of simulation
techniques are in research to develop high-speed simulators for the DSE of
general-purpose computers with less complexity and shorter development time
then conventional cycle accurate simulator. These simulation techniques are
diverse and do not follow one particular pattern. In this section we give
details of some high-level simulation techniques. There might exist other
high-level simulations targeting general-purpose computers.
### 4.1 Interval simulation
Interval simulation [7, 18] is a high-level simulation technique for the DSE
of super-scalar single- and multi- core processors. It raises the level of
abstraction from detailed simulation by using analytical models to derive the
timing simulation of individual cores without the detailed execution of
instructions in all the stages of the pipeline. The model is based on deriving
the execution of an instruction stream in intervals. An interval is decided
based on the miss events e.g. branch misprediction, cache misses, TLB misses
etc. With interval analysis, execution time is partitioned into discrete
intervals using miss events. The analytical models of every core cooperate
with miss events in the system, and can be extended to model the tight
interleaving of threads in multi-core processors.
Interval simulation framework has two parts; functional simulation and timing
simulation, and are connected with each other through a queue. The functional
simulator feeds instructions into the tail of the queue and the timing
simulator reads those instructions from the head of the queue. The functional
simulator generates a dynamic instruction stream, including user-level and
system-level code and is subsequently fed into the timing simulator. The
timing simulator analyzes the code and advances the simulation time as per the
time required to execute an instruction stream. In case of I-cache miss,
branch misprediction and long latency load operations the simulation time is
advanced by the miss latency, branch resolution time plus the front-end
pipeline depth and long latency operations respectively.
#### Discussion
Interval simulator only simulates a small number of cores in super-scalar
machines which disregards hardware microthreading and therefore the complexity
of simulating latency tolerance is not encountered. In the Microgrid we can
have more than 100 cores on the chip, and the architecture is completely
different than super-scalar machine, as it provides fine-grained latency
tolerance based on data-flow scheduling. The way programs can be written for
the Microgrid is also different. Therefore interval simulation can not
directly be used for the DSE of the Microgrid. However, we have learned some
techniques from interval simulation and have used these in HLSim. For
instance, in interval simulation in case of a cache miss the simulation time
is advanced with the addition of cache miss latency. In HLSim we advance the
simulation time with the cache miss latency but adjusted with a latency
tolerance factor based on the number of active threads. Because in case of
latency tolerance the cache miss latency can be shorter than the latency
without any latency tolerance.
### 4.2 Statistical simulation
Statistical simulation has gained interest over the past few years, as it
speeds up simulation by providing short running synthetic traces. The
execution of the original benchmarks is profiled and the key execution
characteristics are captured in a synthetic trace, which closely exhibits
similar execution characteristic as original benchmarks. The key benefit of
statistical simulation is that the synthetic trace clones the dynamic
instruction count with several orders of magnitude smaller than in the
original benchmarks, and therefore reduces the simulation time dramatically.
Nussbaum and Smith [29] and Hughes and Li [19] use statistical simulation
paradigm to evaluate multithreaded programs running on shared-memory
multiprocessor (SMP) systems. They have extended the statistical simulation to
model synchronization and accesses to shared memory. Genbrugge and Eeckhout
[13, 17] use statistical simulation to measure some execution characteristics
in the statistical profile to be able to accurately simulate shared resources
in multi-core processors.
#### Discussion
Statistical simulation is a trace driven simulation technique. A synthetic
trace is generated which can be reduced to a shorter trace and is
representative of the large trace of the benchmarks. The problem with this
technique is that the original trace files can be very large which consume
space and this technique can not consider the dynamic adaptation of multiple
applications on the chip. The high-level simulation of the Microgrid, is
execution driven i.e we dynamically generate events which are representative
of the instruction count in the basic block in a thread. These events are
mapped to the architecture and represent the execution of the application with
fine-grained interleaving. The events have information of a short piece of
code and therefore statistical simulation techniques wee not a suitable choice
to be used in HLSim.
### 4.3 Sampled simulation
The basic idea of sampled simulation is to simulate a number of sampling units
rather than the entire dynamic instruction stream. The sampling units are
selected either randomly [9], periodically [49] or based on phase analysis
[33].
Different research in the multithreaded and multi-core processors simulation
is using sampled simulation. Van Biesbrouck et al. [42] propose the co-phase
matrix for speeding up sampled simultaneous multithreading (SMT) processor
simulation running multi-program workloads. Stenstrom et al. [14] are
researching the premise that fewer sampling units are enough to estimate
overall performance for larger multi-processor systems than for smaller multi-
processor system in case one is interested in aggregate performance only.
Wenisch et al. [48] have obtained similar conclusions of throughput in server
workloads. Barr et al. [3] proposes the Memory Timestamp Record (MTR) to store
micro-architecture state (cache and directory state) at the beginning of the
sampling unit as a checkpoint.
#### Discussion
Sampled simulation is also a trace-based simulation technique which suffers
from the large trace files to be processed and changing an application results
in producing and analysing a different trace file. Every time there is some
optimization in the application, a new trace needs to be generated and
analyzed.
### 4.4 Related works
There are other simulation techniques used in the design space exploration of
general-purpose computers given below.
* •
FPGA prototypes: They have low little simulation time, high accuracy and are
useful in DSE. However these simulations require more development time and are
more complex. They also suffer from combinatoric explosion of considering many
low level parameters during design space exploration. Some examples are: [30,
31, 8, 46].
* •
Trace simulation: These simulation techniques generate large execution traces
from benchmarks, and are used for the evaluation of the architecture. They
avoid the extremely large analysis of the application, by executing the
program only one time, generating the trace and mapping it to the trace to
different configuration of the architecture. However a large storage is
required in order to store the large traces and a change in the application
requires a different trace to be generated. Statistical simulations and
sampled simulations are some of the techniques that addresses the reduction of
the large trace files. Some example are: [9, 20, 27].
## 5 Design space exploration in the Microgrid
### 5.1 Microgrid
The Microgrid [24, 5, 22] is a general-purpose, many-core architecture
developed at the University of Amsterdam which implements hardware multi-
threading using data flow scheduling and a concurrency management protocol in
hardware to create and synchronize threads within and across the cores on
chip. The suggested concurrent programming model for this chip is based on
fork-join constructs, where each created thread can define further concurrency
hierarchically. This model is called the microthreading model and is also
applicable to current multi-core architectures using a library of the
concurrency constructs called _svp-ptl_ [45] built on top of pthreads. In our
work, we focus on a specific implementation of the microthreaded architecture
where each core contains a single issue, in-order RISC pipeline with an ISA
similar to DEC/Alpha, and all cores are connected to an on-chip distributed
memory network [23, 6]. Each core implements the concurrency constructs in its
instruction set and is able to support hundreds of threads and their contexts,
called microthreads and tens of families (i.e. ordered collections of
identical microthreads) simultaneously.
A number of tools and simulators are added to the designer’s toolbox and used
for the evaluation of the Microgrid from different perspective. The compiler
for the Microgrid [25] can generate binary for different implementations of
the Microgrid. We have software libraries that provide the run-time systems
for the microthreading model on the shared memory SMP machines and referred as
_svp-ptl_ [45] and distributed memory for clusters/grids and are referred as
Hydra [28] and _dsvp-ptl_ [44] The SL compiler can generate binary for UTLEON3
[10, 11], MGSim [6, 32] and HLSim [37, 38, 39, 36, 40, 41, 34, 35].
HLSim is a high-level simulation technique aimed for the DSE of the Microgrid
and is based on discrete event simulation technique. It is execution driven
simulator and therefore does not suffer from the large size of trace files.
The events are dynamically mapped to the architecture at run time. We have
built the simulator from scratch without using any off-the-shelf code, but
some simulation techniques from Sesame and Interval simulation were used
during the development for inspiration.
### 5.2 Design space in the Microgrid
The Microgrid is a complex many-cores architecture and therefore has a huge
design space for complex application. In order to have an efficient and
validated system in the silicon we need to perform DSE in the Microgrid to
explore the performance of different applications on the different
configurations of the architecture. After DSE we can perform design pruning to
change these parameters that affect the performance. We categorize the design
space of the Microgrid as:
* •
Static architectural parameters:
1. 1.
Thread table size
2. 2.
Family table size
3. 3.
Frequency of cores and memory
4. 4.
Number of cores sharing an FPU
5. 5.
Frequency of delegation and distribution network
6. 6.
Size of L1-cache and L2-cache
7. 7.
Associativity of L1-cache and L2-cache
8. 8.
Number of L1-caches sharing L2-cache
9. 9.
Number of L2-caches in low-level ring
10. 10.
Number of low-level rings associated in the top-level ring
11. 11.
Distribution of address space of RAM into banks
12. 12.
Size of directory and root directory
13. 13.
The memory architecture
14. 14.
Synchronization-aware protocol
* •
Dynamic application parameters:
1. 1.
Place size
2. 2.
Window size
3. 3.
Cold caches
There are some other parameters that are very low-level e.g. size of the chip,
FPU frequency, pipeline stages, the way cores are distributed on the chip etc.
We have shown only the parameters that we will simulate in the current
implementation of HLSim for the design space exploration in the Microgrid.
## 6 Conclusion
DSE is required in all kind of computer systems. The use of high-level
simulators for DSE is pioneered in embedded systems and getting popular in
general-purpose systems. As the Microgrid has a huge design space therefore,
low-level simulators are not justifiable to be used for design space
exploration. We need high-level simulators for the efficient design space
exploration.
## Acknowledgement
The author would like to thank Dr. Raphael Poss, Dr. Michiel van Tol and Prof.
dr. Chris Jesshope.
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|
arxiv-papers
| 2013-09-22T02:01:28 |
2024-09-04T02:49:51.299813
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Irfan Uddin",
"submitter": "Irfan Uddin",
"url": "https://arxiv.org/abs/1309.5551"
}
|
1309.5589
|
# A generalization of Ćirić fixed point theorems
Nguyen Van Dung Faculty of Mathematics and Information Technology Teacher
Education, Dong Thap University, Cao Lanh, Dong Thap 871200, Viet Nam
[email protected], [email protected] , Poom Kumam Department of
Mathematics, Faculty of Science, King Mongkut’s University of Technology
Thonburi (KMUTT), Bang Mod, Thrung Khru, Bangkok 10140, Thailand
[email protected] and Kanokwan Sitthithakerngkiet Department of
Mathematics, Faculty of Applied Science, King Mongkut’s University of North
Bangkok (KMUTNB), Wongsawang, Bangsue, Bangkok 10800, Thailand
[email protected]
###### Abstract.
In this paper, we state and prove a generalization of Ćirić fixed point
theorems in metric space by using a new generalized quasi-contractive map.
These theorems extend other well known fundamental metrical fixed point
theorems in the literature (Banach [1], Kannan [11], Nadler [13], Reich [15],
etc.) Moreover, a multi-valued version for generalized quasi-contraction is
also established.
###### Key words and phrases:
Ćirić fixed point, metric space
###### 2000 Mathematics Subject Classification:
Primary 47H10, 54H25; Secondary 54D99, 54E99
## 1\. Introduction
The Banach’s contraction principle [1] which was first appeared in 1922 is one
of the most useful and important theorems in classical functional analysis.
Its utility is not only to prove that, in a complete metric space $X$, the
contraction map $T$ (i.e., $d(Tx,Ty)\leq\alpha d(x,y)$ for some
$0\leq\alpha<1$ and for all $x,y\in X$) has a unique fixed point but also to
show that the Picard iteration converges to the fixed point. For the reason
that the contraction must be continuous, there are many researchers establish
the fixed point theorems on various classes of operators that are weaker than
contractive conditions but are not continuous, see for example [11, 15].
One of the most well-known results in generalizations of Banach’s contraction
principle which the Picard iteration still converges to the fixed point of map
is the Ćirić fixed point theorem [4]. Before providing the Ćirić fixed point
theorem, we recall that a self-map $T$ on a metric space $(X,d)$, is said to
be a _quasi-contraction_ iff there exists a nonnegative number $q<1$ such that
for all $x,y\in X$,
(1.1) $\displaystyle d(Tx,Ty)$ $\displaystyle\leq$ $\displaystyle
q\max\big{\\{}d(x,y),d(x,Tx),d(y,Ty),d(x,Ty),d(y,Tx)\big{\\}}.$
The Ćirić fixed point theorem is given by the following theorem.
###### Theorem 1.1 ([4], Theorem 1).
Let the metric space $X$ be $T$-orbitally complete and let $T$ be a quasi-
contraction. Then we have
1. (1)
$T$ has a unique fixed point $x^{*}$ in $X$.
2. (2)
$\lim\limits_{n\rightarrow\infty}T^{n}x=x^{*}$ for all $x\in X$.
3. (3)
$d(T^{n}x,x^{*})\leq\cfrac{q^{n}}{1-q}d(x,Tx)$ for all $x\in X$.
This result was generalized to many results, such as a common fixed point
theorem of nonlinear contraction [17, Theorem 4], a generalized
$\varphi$-contraction [2, Section 2.6], a Ćirić almost contraction [3, Theorem
3.2] and see also [5, 10, 12]. But from the well-known result of Rhoades [14]
in 1977 to recent surveys, in Berinde [2] and Collaco and Silva [6] for
instance, there were no any other value added to quasi-contraction condition.
On the other hand, the Banach’s contraction principle has been extended to
multi-valued contractions by Nadler [13] and see also [7, 8, 9, 16].
In this paper, we define a new generalized quasi-contraction by adding four
new values $d(T^{2}x,x)$, $d(T^{2}x,Tx)$, $d(T^{2}x,y)$, $d(T^{2}x,Ty)$ to a
quasi-contraction condition. Also, an example is presented. After that we
state and prove unique fixed point theorems which are the generalization of
Ćirić fixed point theorem in [4]. Moreover, we also establish fixed point
theorems for multi-valued generalized quasi-contraction.
## 2\. Preliminaries
First, we recall some notions which will be used in what follows. Let $(X,d)$
be a metric space and $A,B$ be any two subsets of $X$. We denote
$\displaystyle D(A,B)$ $\displaystyle=$ $\displaystyle\inf\big{\\{}d(a,b):a\in
A,b\in B\big{\\}}$ $\displaystyle\rho(A,B)$ $\displaystyle=$
$\displaystyle\sup\big{\\{}d(a,b):a\in A,b\in B\big{\\}}$ $\displaystyle
BN(X)$ $\displaystyle=$ $\displaystyle\big{\\{}A:\emptyset\neq A\subset
X\text{ and }\delta(A)<+\infty\big{\\}},$
where $\delta(A):=\sup\big{\\{}d(a,b):a,b\in A\big{\\}}$.
###### Definition 2.1 ([4]).
Let $T:X\longrightarrow X$ be a map on metric space. For each $x\in X$ and for
any positive integer $n$, put
$O_{T}(x,n)=\\{x,Tx,\ldots,T^{n}x\\}~{}and~{}O_{T}(x,+\infty)=\\{x,Tx,\ldots,T^{n}x,\ldots\\}.$
The set $O_{T}(x,+\infty)$ is called the _orbit_ of $T$ at $x$ and the metric
space $X$ is called _$T$ -orbitally complete_ if every Cauchy sequence in
$O_{T}(x,+\infty)$ is convergent in $X$.
Note that every complete metric space is $T$-orbitally complete for all maps
$T:X\longrightarrow X$. The following example shows that there exists a
$T$-orbitally complete metric space but it is not complete.
###### Example 2.2.
Let $(X,d)$ be a metric space which is not complete and $T:X\longrightarrow X$
be the map defined by $Tx=x_{0}$ for all $x\in X$ and some $x_{0}\in X$. Then
$(X,d)$ is a $T$-orbitally complete metric space which is not complete.
###### Definition 2.3 ([4]).
Let $F:X\longrightarrow BN(X)$ be a multi-valued mapping. Let $x_{0}\in X$, an
_orbit_ of $F$ at $x_{0}$ is a sequence
$\big{\\{}x_{n}:x_{n}\in Fx_{n-1},n\in\mathbb{N}\big{\\}}.$
A space $X$ is called to be _$F$ -orbitally complete_ if every Cauchy sequence
which is a subsequence of an orbit of $F$ at $x$ for some $x\in X$, converges
in $X$.
Next, the definitions of generalized quasi-contraction for single-valued and
multi-valued are given as follows;
###### Definition 2.4.
Let $T:X\longrightarrow X$ be a mapping on metric space $X$. The mapping $T$
is said to be a _generalized quasi-contraction_ iff there exists $q\in[0,1)$
such that for all $x,y\in X$,
$\displaystyle d(Tx,Ty)$ $\displaystyle\leq$ $\displaystyle
q\max\big{\\{}d(x,y),d(x,Tx),d(y,Ty),d(x,Ty),d(y,Tx),$ $\displaystyle
d(T^{2}x,x),d(T^{2}x,Tx),d(T^{2}x,y),d(T^{2}x,Ty)\big{\\}}.$
###### Example 2.5.
Let $X=\\{1,2,3,4,5\\}$ with $d$ defined as
$d(x,y)=\left\\{\begin{array}[]{ll}0&\hbox{ if }x=y\\\ 2&\hbox{ if
}(x,y)\in\big{\\{}(1,4),(1,5),(4,1),(5,1)\big{\\}}\\\ 1&\hbox{
otherwise.}\end{array}\right.$
Let $T:X\longrightarrow X$ be defined by
$T1=T2=T3=1,T4=2,T5=3.$
Then, we have
$d(Tx,Ty)=d(1,1)=0~{}\text{ if }~{}x,y\in\\{1,2,3\\};$
$d(T1,T4)=d(T2,T4)=d(T3,T4)=d(1,2)=1;$ $d(T1,4)=d(T2,4)=d(T3,4)=d(1,4)=2;$
$d(T1,T5)=d(T2,T5)=d(T3,T5)=d(1,3)=1;$ $d(T1,5)=d(T2,5)=d(T3,5)=d(1,5)=2;$
$d(T4,T5)=d(2,3)=1;$ $d(4,5)=d(4,T4)=d(5,T5)=d(4,T5)=d(5,T4)=1;$
$d(T^{2}4,4)=d(T2,4)=d(1,4)=2;$ $d(T^{2}5,5)=d(T3,5)=d(1,5)=2.$
The above calculations show that $T$ is not quasi-contraction for $x=4$ and
$y=5$ because there is no a nonnegative number $q<1$ satisfying the equation
(2.4). However, $T$ is generalized quasi-contraction since the (2.4) holds for
some $q\in\big{[}0.5,1\big{)}$ and for all $x,y\in X$.
## 3\. The main results
On the following results, we state and prove the new fixed point theorems
which are general cases of the Ćirić fixed point theorem.
###### Theorem 3.1.
Let $(X,d)$ be a metric space. Suppose that $T:X\longrightarrow X$ is a
generalized quasi-contraction and $X$ is $T$-orbitally complete. Then we have
1. (1)
$T$ has a unique fixed point $x^{*}$ in $X$.
2. (2)
$\lim\limits_{n\rightarrow\infty}T^{n}x=x^{*}$ for all $x\in X$.
3. (3)
$d(T^{n}x,x^{*})\leq\cfrac{q^{n}}{1-q}d(x,Tx)$ for all $x\in X$ and
$n\in\mathbb{N}$.
###### Proof.
(1). Step 1. _$T$ has a fixed point._ For each $x\in X$ and $1\leq i\leq n-1$
and $1\leq j\leq n$, we have
$\displaystyle d(T^{i}x,T^{j}x)$ $\displaystyle=$ $\displaystyle
d(TT^{i-1}x,TT^{j-1}x)$ $\displaystyle\leq$ $\displaystyle
q\max\big{\\{}d(T^{i-1}x,T^{j-1}x),d(T^{i-1}x,TT^{i-1}x),d(T^{j-1}x,TT^{j-1}x),d(T^{i-1}x,TT^{j-1}x),$
$\displaystyle
d(T^{j-1}x,TT^{i-1}x),d(T^{2}T^{i-1}x,T^{i-1}x),d(T^{2}T^{i-1}x,TT^{i-1}x),d(T^{2}T^{i-1}x,T^{j-1}x)$
$\displaystyle d(T^{2}T^{i-1}x,TT^{j-1}x)\big{\\}}$ $\displaystyle=$
$\displaystyle
q\max\big{\\{}d(T^{i-1}x,T^{j-1}x),d(T^{i-1}x,T^{i}x),d(T^{j-1}x,T^{j}x),d(T^{i-1}x,T^{j}x),$
$\displaystyle
d(T^{j-1}x,T^{i}x),d(T^{i+1}x,T^{i-1}x),d(T^{i+1}x,T^{i}x),d(T^{i+1}x,T^{j-1}x)$
$\displaystyle d(T^{i+1}x,T^{j}x)\big{\\}}$ $\displaystyle\leq$ $\displaystyle
q\delta\big{[}O_{T}(x,n)\big{]}$
where $\delta\big{[}O_{T}(x,n)\big{]}=\max\big{\\{}d(T^{i}x,T^{j}x):0\leq
i,j\leq n\big{\\}}$.
From (3), since $0\leq q<1$, there exists $k_{n}(x)\leq n$ such that
(3.2) $d(x,T^{k_{n}(x)}x)=\delta\big{[}O_{T}(x,n)\big{]}.$
Then we have
$\displaystyle d(x,T^{k_{n}(x)}x)$ $\displaystyle\leq$ $\displaystyle
d(x,Tx)+d(Tx,T^{k_{n}(x)}x)$ $\displaystyle\leq$ $\displaystyle
d(x,Tx)+q\delta\big{[}O_{T}(x,n)\big{]}$ $\displaystyle=$ $\displaystyle
d(x,Tx)+qd(x,T^{k_{n}(x)}x).$
It implies that
(3.3)
$\delta\big{[}O_{T}(x,n)\big{]}=d(x,T^{k_{n}(x)}x)\leq\cfrac{1}{1-q}d(x,Tx).$
For all $n,m\leq 1$ and $n<m$, it follows from the generalized quasi-
contractive condition of $T$ and (3.3) that
$\displaystyle d(T^{n}x,T^{m}x)$ $\displaystyle=$ $\displaystyle
d(TT^{n-1}x,T^{m-n+1}T^{n-1}x)$ $\displaystyle\leq$ $\displaystyle
q\delta\big{[}O_{T}(T^{n-1}x,m-n+1)\big{]}$ $\displaystyle=$ $\displaystyle
qd(T^{n-1}x,T^{k_{m-n+1}(T^{n-1}x)}T^{n-1}x)$ $\displaystyle=$ $\displaystyle
qd(TT^{n-2}x,T^{k_{m-n+1}(T^{n-1}x)+1}T^{n-2}x)$ $\displaystyle\leq$
$\displaystyle q^{2}\delta\big{[}O_{T}(T^{n-2}x,k_{m-n+1}(T^{n-1}x)+1)\big{]}$
$\displaystyle\leq$ $\displaystyle
q^{2}\delta\big{[}O_{T}(T^{n-2}x,m-n+2)\big{]}$ $\displaystyle\leq$
$\displaystyle\ldots$ $\displaystyle\leq$ $\displaystyle
q^{n}\delta\big{[}O_{T}(x,m)\big{]}$ $\displaystyle\leq$
$\displaystyle\cfrac{q^{n}}{1-q}d(x,Tx).$
Since $\lim\limits_{n\rightarrow\infty}q^{n}=0$, $\\{T^{n}x\\}$ is a Cauchy
sequence in $X$. Since $X$ is $T$-orbitally complete, there exists $x^{*}\in
X$ such that
(3.5) $\lim\limits_{n\rightarrow\infty}T^{n}x=x^{*}.$
By using the generalized quasi-contractive condition of $T$ again, we have
$\displaystyle d(x^{*},Tx^{*})$ $\displaystyle\leq$ $\displaystyle
d(x^{*},T^{n+1}x)+d(T^{n+1}x,Tx^{*})$ $\displaystyle=$ $\displaystyle
d(x^{*},T^{n+1}x)+d(TT^{n}x,Tx^{*})$ $\displaystyle\leq$ $\displaystyle
d(x^{*},T^{n+1}x)+q\max\big{\\{}d(T^{n}x,x^{*}),d(T^{n}x,TT^{n}x),d(x^{*},Tx^{*}),d(T^{n}x,Tx^{*}),$
$\displaystyle
d(x^{*},TT^{n}x),d(T^{2}T^{n}x,T^{n}x),d(T^{2}T^{n}x,TT^{n}x),d(T^{2}T^{n}x,x^{*}),d(T^{2}T^{n}x,Tx^{*})\big{\\}}$
$\displaystyle=$ $\displaystyle
d(x^{*},T^{n+1}x)+q\max\big{\\{}d(T^{n}x,x^{*}),d(T^{n}x,T^{n+1}x),d(x^{*},Tx^{*}),d(T^{n}x,Tx^{*}),$
$\displaystyle
d(x^{*},T^{n+1}x),d(T^{n+2}x,T^{n}x),d(T^{n+2}x,T^{n+1}x),d(T^{n+2}x,x^{*}),d(T^{n+2}x,Tx^{*})\big{\\}}.$
Taking the limit as $n\rightarrow\infty$ in (3), and using (3.5), we get
$d(x^{*},Tx^{*})\leq qd(x^{*},Tx^{*}).$ Since $q\in[0,1)$, we obtain
$d(x^{*},Tx^{*})=0$, that is, $x^{*}=Tx^{*}$. Then $T$ has a fixed point.
Step 2. _The fixed point of $T$ is unique._ Let $x^{*},y^{*}$ be two fixed
points of $T$. Since $T$ is generalized quasi-contraction, we have
$\displaystyle d(x^{*},y^{*})$ $\displaystyle=$ $\displaystyle
d(Tx^{*},Ty^{*})$ $\displaystyle\leq$ $\displaystyle
q\max\big{\\{}d(x^{*},y^{*}),d(x^{*},Tx^{*}),d(y^{*},Ty^{*}),d(x^{*},Ty^{*}),d(y^{*},Tx^{*}),$
$\displaystyle
d(T^{2}x^{*},x^{*}),d(T^{2}x^{*},Tx^{*}),d(T^{2}x^{*},y^{*}),d(T^{2}x^{*},Ty^{*})\big{\\}}$
$\displaystyle=$ $\displaystyle qd(x^{*},y^{*}).$
Since $q\in[0,1)$, we obtain $d(x^{*},y^{*})=0$. That is, $x^{*}=y^{*}$. Then
the fixed point of $T$ is unique.
(2). It is proved by (3.5).
(3). Taking the limit as $m\rightarrow\infty$ in (3), we get
$d(T^{n}x,x^{*})\leq\cfrac{q^{n}}{1-q}d(x,Tx).$ ∎
###### Corollary 3.2.
Let $(X,d)$ be a metric space and $T:X\longrightarrow X$ be a map satisfying
the following:
1. (1)
$X$ is $T$-orbitally complete.
2. (2)
There exists $k\in\mathbb{N}$ and $q\in[0,1)$ such that for all $x,y\in X$,
$\displaystyle d(T^{k}x,T^{k}y)$ $\displaystyle\leq$ $\displaystyle
q\max\big{\\{}d(x,y),d(x,T^{k}x),d(y,T^{k}y),d(x,T^{k}y),d(y,T^{k}x),$
$\displaystyle
d(T^{2k}x,x),d(T^{2k}x,T^{k}x),d(T^{2k}x,y),d(T^{2k}x,T^{k}y)\big{\\}}.$
Then we have
1. (1)
$T$ has a unique fixed point $x^{*}$ in $X$.
2. (2)
$d(T^{n}x,x^{*})\leq\cfrac{q^{m}}{1-q}\max\big{\\{}d(T^{i}x,T^{i+k}x):i=0,1,\ldots,k-1\big{\\}}$
for all $x\in X$ and $n\in\mathbb{N}$ where $m$ is the greatest integer not
exceeding $\cfrac{n}{k}$.
3. (3)
$\lim\limits_{n\rightarrow\infty}T^{n}x=x^{*}$ for all $x\in X$.
###### Proof.
(1). By the conclusion of Theorem 3.1, $T^{k}$ has a unique fixed point
$x^{*}$ and $T^{k}(Tx^{*})=T(T^{k}x^{*})=Tx^{*}$. It implies that
$Tx^{*}=x^{*}$, that is, $T$ has a fixed point $x^{*}$. The uniqueness of the
fixed point of $T$ is easy to see.
(2). Let $n\in\mathbb{N}$. Then $n=mk+j$, $0\leq j<k$ and for each $x\in X$,
$T^{n}x=(T^{k})^{m}T^{j}x$. It follows from Theorem 3.1.(3) that
$\displaystyle d(T^{n}x,x^{*})$ $\displaystyle\leq$
$\displaystyle\cfrac{q^{m}}{1-q}d(T^{j}x,T^{k}T^{j}x)$ $\displaystyle\leq$
$\displaystyle\cfrac{q^{m}}{1-q}\max\big{\\{}d(T^{i}x,T^{i+k}x):i=0,1,\ldots,k-1\big{\\}}.$
(3). It is a direct consequence of (2). ∎
###### Corollary 3.3 ([4], Theorem 2).
Let $(X,d)$ be a metric space and $T:X\longrightarrow X$ be a map satisfying
the following:
1. (1)
$X$ is $T$-orbitally complete.
2. (2)
There exists $k\in\mathbb{N}$ and $q\in[0,1)$ such that for all $x,y\in X$,
(3.8) $\displaystyle d(T^{k}x,T^{k}y)$ $\displaystyle\leq$ $\displaystyle
q\max\big{\\{}d(x,y),d(x,T^{k}x),d(y,T^{k}y),d(x,T^{k}y),d(y,T^{k}x)\big{\\}}.$
Then we have
1. (1)
$T$ has a unique fixed point $x^{*}$ in $X$;
2. (2)
$d(T^{n}x,x^{*})\leq\cfrac{q^{m}}{1-q}\max\big{\\{}d(T^{i}x,T^{i+k}x):i=0,1,\ldots,k-1\big{\\}}$
for all $x\in X$ and $n\in\mathbb{N}$ where $m$ is the greatest integer not
exceeding $\cfrac{n}{k}$;
3. (3)
$\lim\limits_{n\rightarrow\infty}T^{n}x=x^{*}$ for all $x\in X$.
Now, we denote the multi-valued mapping $F:X\longrightarrow BN(X)$ of
generalized quasi-contraction by
$\displaystyle\rho(Fx,Fy)$ $\displaystyle\leq$ $\displaystyle
q\max\big{\\{}d(x,y),\rho(x,Fx),\rho(y,Fy),D(x,Fy),D(y,Fx),$ $\displaystyle
D(F^{2}x,x),D(F^{2}x,Fx),D(F^{2}x,y),D(F^{2}x,Fy)\big{\\}},$
for some $q\in[0,1)$ and for all $x,y\in X$. The following theorem presents
the fixed point theorem for multi-valued version of generalized quasi-
contractive mapping.
###### Theorem 3.4.
Let $(X,d)$ be a metric space and $F:X\longrightarrow BN(X)$ be a multi-valued
map. Suppose that $F$ is a generalized quasi-contraction and $X$ is
$F$-orbitally complete. Then we have
1. (1)
$F$ has a unique fixed point $x^{*}$ in $X$ and $Fx^{*}=\\{x^{*}\\}$.
2. (2)
For each $x_{0}\in X$, there exists an orbit $\\{x_{n}\\}_{n}$ of $F$ at
$x_{0}$ such that $\lim\limits_{n\rightarrow\infty}x_{n}=x^{*}$ for all $x\in
X$, and
3. (3)
$d(x_{n},x^{*})\leq\cfrac{(q^{1-a})^{n}}{1-q^{1-a}}d(x_{0},x_{1})$ for all
$n\in\mathbb{N}$, where $a<1$ is any fixed positive number.
###### Proof.
(1). Given $a\in(0,1)$ and defined a single-valued mapping $T:X\longrightarrow
X$ by the following statement:
$for~{}each~{}x\in X,~{}Tx\in Fx~{}~{}satisfies~{}~{}d(x,Tx)\geq
q^{a}\rho(x,Fx).$
By the Definition 2.3 and the condition of $F$, we have for every $x,y\in X$,
$\displaystyle d(Tx,Ty)$ $\displaystyle\leq$ $\displaystyle\rho(Fx,Fy)$
$\displaystyle\leq$ $\displaystyle
q\max\big{\\{}d(x,y),\rho(x,Fx),\rho(y,Fy),D(x,Fy),D(y,Fx),$ $\displaystyle
D(F^{2}x,x),D(F^{2}x,Fx),D(F^{2}x,y),D(F^{2}x,Fy)\big{\\}}$ $\displaystyle=$
$\displaystyle
qq^{-a}\max\big{\\{}q^{a}d(x,y),q^{a}\rho(x,Fx),q^{a}\rho(y,Fy),q^{a}D(x,Fy),q^{a}D(y,Fx),$
$\displaystyle
q^{a}D(F^{2}x,x),q^{a}D(F^{2}x,Fx),q^{a}D(F^{2}x,y),q^{a}D(F^{2}x,Fy)\big{\\}}$
$\displaystyle\leq$ $\displaystyle
q^{1-a}\max\big{\\{}d(x,y),d(x,Tx),d(y,Ty),d(x,Ty),d(y,Tx),$ $\displaystyle
d(T^{2}x,x),d(T^{2}x,Tx),d(T^{2}x,y),d(T^{2}x,Ty)\big{\\}}.$
By Theorem 3.1, we conclude that $T$ has a unique fixed point $x^{*}$. Then
$\rho(x^{*},Fx^{*})\leq q^{a}d(x^{*},Tx^{*})=0$ implies that
$\rho(x^{*},Fx^{*})=0$. Then $x^{*}$ is a fixed point of $F$ and
$Fx^{*}=\\{x^{*}\\}$. From the direct consequences of Theorem 3.1 where
$x_{n}=T^{n}x$ for all $n\in\mathbb{N}$, we obtain that (2) and (3) hold. ∎
###### Corollary 3.5 ([4], Theorem 3).
Let $(X,d)$ be a metric space and $F:X\longrightarrow BN(X)$ be a multi-valued
map satisfying the following:
1. (1)
$X$ is $F$-orbitally complete.
2. (2)
There exists $q\in[0,1)$ such that for all $x,y\in X$,
(3.10) $\displaystyle\rho(Fx,Fy)$ $\displaystyle\leq$ $\displaystyle
q\max\big{\\{}d(x,y),\rho(x,Fx),\rho(y,Fy),D(x,Fy),D(y,Fx)\big{\\}}.$
Then we have
1. (1)
$F$ has a unique fixed point $x^{*}$ in $X$ and $Fx^{*}=\\{x^{*}\\}$.
2. (2)
For each $x_{0}\in X$, there exists an orbit $\\{x_{n}\\}_{n}$ of $F$ at
$x_{0}$ such that $\lim\limits_{n\rightarrow\infty}x_{n}=x^{*}$ for all $x\in
X$, and
3. (3)
$d(x_{n},x^{*})\leq\cfrac{(q^{1-a})^{n}}{1-q^{1-a}}d(x_{0},x_{1})$ for all
$n\in\mathbb{N}$, where $a<1$ is any fixed positive number.
###### Example 3.6.
Let $(X,d)$ and $T:X\longrightarrow X$ be defined by Example 2.5.
It is easy to see that $X$ is $T$-orbitally complete metric space. By the
definition of the distance $d$ and mapping $T$ , we conclude that $X$ and $T$
satisfy all of the conditions in Theorem 3.1. Clearly, $x^{*}=1$ is a unique
fixed point of $T$.
Note that, if $x\in\\{1,2,3\\}$ then $T^{n}x=1$ for $n=1,2,3,...$ and if
$x\in\\{4,5\\}$ then $T^{n}x=1$ for $n=2,3,4,...$. That is
$\lim\limits_{n\rightarrow\infty}T^{n}x=x^{*}$ for all $x\in X$.
Let $q\in\big{[}0.5,1\big{)}$ be fixed by the generalized quasi-contraction of
$T$ which arises from Example 2.5. We see that the inequality
$d(T^{n}x,x^{*})\leq\cfrac{q^{n}}{1-q}d(x,Tx)$ holds for all $x\in X$ and
$n\in\mathbb{N}$.
Therefore, this example is presented to certify the results of Theorem 3.1.
However, it is not applicable to Theorem 1.1.
###### Remark 3.7.
Example 3.6 shows that our results are proper generalizations of Ćirić fixed
point theorems in [4]. Then our results are exactly a new form of fixed point
theorems in metric spaces. Moreover, we may generalize other fixed point
theorems contained at most five mentioned values in the literature to that
contain $d(T^{2}x,x)$, $d(T^{2},Tx)$, $d(T^{2}x,y)$, $d(T^{2}x,Ty)$ in
addition.
Acknowledgement. This research of is supported by King Mongkut’s University of
Technology North Bangkok, Thailand
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|
arxiv-papers
| 2013-09-22T10:53:43 |
2024-09-04T02:49:51.311516
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Nguyen Van Dung and Poom Kumam and Kanokwan Sitthithakerngkiet",
"submitter": "Dung Nguyen Van",
"url": "https://arxiv.org/abs/1309.5589"
}
|
1309.5649
|
# Project Description
Gunnar Carlsson 111Research supported in part by NSF DMS-0406992
Department of Mathematics, Stanford University
Stanford, California 94305
# A geometric reformulation of the representational assembly conjecture for
the $K$-theory of fields with pro-$l$ absolute Galois group
Gunnar Carlsson 111Research supported in part by NSF DMS-0406992
Department of Mathematics, Stanford University
Stanford, California 94305
## 1 Introduction
In recent years there have been striking developments in the study of the
algebraic $K$-theory of fields (see [41], [16], [9]). These developments have
resulted in the identification of the $E_{2}$-term of a spectral sequence
converging to the algebraic $K$-theory of fields (rather a mod-$p$ version of
it) in terms of Galois cohomology, and have resolved long standing conjectures
of Milnor and Bloch-Kato in the area. In particular the $E_{2}$-term depends
only on the absolute Galois group of the field. These results raise a pair of
questions.
* •
For a field $F$, is there a construction depending only on the absolute Galois
group $G_{F}$ of $F$, which reconstructs the full $K$-theory spectrum of $F$,
rather than an $E_{2}$-term of a spectral sequence? Can such a construction be
made functorial for inclusions of subfields and the corresponding inclusions
of absolute Galois groups? Is the construction readily describable in terms of
well understood constructions on the absolute Galois group?
* •
The conjectures of Milnor and Bloch-Kato are cohomological statements about
the absolute Galois group, specifically that the cohomology is generated in
degree one and that the relations are quadratic. Are there properties more
closely connected to the structure of the absolute Galois group as a group
which are relevant?
A conjectural answer to the first question in the case of geometric fields
(i.e. fields containing an algebraically closed subfield $k$) was proposed in
[8], and proved in the case of abelian absolute Galois groups. This conjecture
is formulated as follows. We assume we are given a field $F$ containing an
algebraically closed subfield $k\subseteq F$, let $\overline{F}$ denote its
algebraic closure, and let $G_{F}$ denote its absolute Galois group
$Gal(\overline{F}/F)$, in general a profinite group. One can construct a
version of the equivariant $K$-theory spectrum, $K^{G}$, which applies to
profinite groups, and observe that $K^{G}(\overline{F})\cong K(F)$. On may
also consider the natural map $K^{G}(k)\rightarrow K^{G}(\overline{F})$ of
spectra. This map is far from being an equivalence since, for example,
$\pi_{0}K^{G}(k)$ is the free abelian group on the continuous $k$-linear
representations of $G_{F}$, where $k$ is given the discrete topology, and on
the other hand
$\pi_{0}K^{G}(\overline{F})\cong\pi_{0}K(F)\cong\mathbb{Z}$
However, the spectra $K^{G}(k)$ and $K^{G}(\overline{F})$ both carry
multiplicative structures, which makes $K^{G}(\overline{F})$ into a module
over the commutative ring spectrum $K^{G}(k)$. See [15] and [23] for a
discussion of these multiplicative theories. In [7], a derived version of the
completion construction for commutative rings was developed, which associates
to a homomorphism of commutative ring spectra $f:A\rightarrow B$ and an
$A$-module $M$ a completion of $M$ along the homomorphism $f$, denoted by
$M^{\wedge}_{f}$, which agrees with the standard completion construction for
finitely generated modules over Noetherian commutative rings. We will denote
by $\varepsilon$ the homomorphism $K^{G}(k)\rightarrow\mathbb{H}_{l}$ which is
induced by the functor which assigns to every representation of $G_{F}$ its
dimension mod $l$, where $\mathbb{H}_{l}$ is the mod-$l$ Eilenberg-MacLane
spectrum, regarded as the result of an infinite loop space machine applied to
the symmetric monoidal category with object set equal to
$\mathbb{Z}/l\mathbb{Z}$, and with only identity morphisms. At this point, we
may construct derived completions along the homomorphism $\varepsilon$, to
obtain a map
$\alpha_{F}:K^{G}(k)^{\wedge}_{\varepsilon}\rightarrow
K^{G}(\overline{F})^{\wedge}_{\varepsilon}$
called the representational assembly, and the conjecture is that this map is
an equivalence of spectra.
The difficulty in approaching this conjecture is that we do not have simple
direct methods for computing with the completions. We will instead approach
the problem indirectly as follows. We will construct an affine scheme ${\cal
E}{G_{F}}$, with associated ring equipped with the discrete topology, and
equipped with a continuous action by the topological group $G_{F}$.
Equivariant $K$-theory can be suitably defined for such actions. The scheme
${\cal E}{G_{F}}$ has the following properties.
1. 1.
The maps
$K^{G_{F}}(Spec(k))^{\wedge}_{\varepsilon}\rightarrow K^{G_{F}}({\cal
E}{G_{F}})^{\wedge}_{\varepsilon}$
and
$K^{G_{F}}(Spec(\overline{F}))^{\wedge}_{\varepsilon}\rightarrow
K^{G_{F}}(Spec(\overline{F})\times{\cal E}{G_{F}})^{\wedge}_{\varepsilon}$
are equivalences.
2. 2.
There is a well-defined notion orbit scheme for ${\cal E}F$, and there are
equivalences
$K({\cal E}{G_{F}}/G_{F})\rightarrow K^{G_{F}}({\cal E}{G_{F}})$
and
$K(Spec(\overline{F})\mathrel{\mathop{\times}\limits_{G_{F}}}{\cal
E}{G_{F}})\rightarrow K^{G_{F}}(Spec(\overline{F})\times{\cal E}{G_{F}})$
This follows from a freeness property of ${\cal E}{G_{F}}$.
3. 3.
There is a natural equivalence $K^{G_{F}}({{\cal
E}G_{F})}^{\wedge}_{\epsilon}\cong K^{G_{F}}({\cal E}G_{F})^{\wedge}_{l}$,
where $(-)^{\wedge}_{l}$ denotes Bousfied-Kan $l$-adic completion.
These properties give rise to the following diagram.
$\begin{diagram}$
The goal is to prove that $\alpha_{F}$ is an equivalence. Properties (1) and
(3) above show that all the horizontal maps are equivalences, and it therefore
follows that it will suffice to prove that the right hand vertical map $\eta$
is an equivalence. Proving that the map $\eta$ is an equivalence can be
thought of as a parametrized rigidity theorem. We will use this reduction in
the paper [10] to prove the conjecture. Some comments on the construction are
in order.
* •
${\cal E}G$ is not functorial for homomorphisms of groups, and it depends on a
number of choices. Moreover, ${\cal E}G_{F}$ is not even
$\mathbb{A}^{1}$-contractible, as its algebraic $K$-theory is not isomorphic
to the algebraic $K$-theory of the base field $k$. What can be proved is that
it is potentially $l$-adically contractible, in that $K({\cal E}G_{F})$ has
$l$-adic completion equivalent to the $l$-adic completion of $K(k)$. We
believe that there should ultimately be canonical constructions for the
universal space of a profinite group, and that our construction is a
particular model which is $l$-adically equivalent to it.
* •
The equivariant $K$-theory of ${\cal E}G_{F}$ is interesting in that the
natural action of the representation ring of $G_{F}$, $R[G_{F}]$, on
$K_{*}({\cal E}G_{F})$ has the property that the augmentation ring of
$R[G_{F}]$ acts trivially on it. This fact is crucial in the proof of Property
(3) above.
* •
The construction of ${\cal E}{G_{F}}$ relies on a structural property shared
by all absolute Galois groups of geometric fields, that of total torsion
freeness, which asserts that the abelianization of any closed subgroup of
$G_{F}$ is torsion free. This is a surprisingly restrictive property for
groups, and permits the construction of free actions on varieties of the form
$\mathbb{G}_{m}^{n}$, which turn out to be extremely useful.
* •
${\cal E}G_{F}$ is of the form $Spec(A)$, where
$A=\bigcup_{n}k[t_{s}^{\pm\frac{1}{l^{n}}};s\in S]$, where $S$ is some
indexing set.
* •
There has been a great deal of work on reconstructing the arithmetic of a
fields from its Galois group, see [30] and [4], for example. In the abstract,
this asserts that one is able to reconstruct the $K$-theory spectrum of the
field from its absolute Galois group. The point of the present work is to
identify the $K$-theory spectrum in terms of well understood constructions on
the absolute Galois group, with the expectation that such understanding will
make applying $K$-theory to arithmetic problems will be made easier.
* •
There is another formulation of the main conjecture, which would replace the
ring spectra which occur by their counterparts within equivariant spectra, and
would therefore naturally replace all homotopy groups which occur by their
corresponding Mackey functors, and further replace the representation rings by
their Green functor counterparts (see [5]). The two conjectures are not a
priori equivalent, since the relationship between the completion in the
category of rings and the completion in the Green functor context is not at
all clear. Indeed, our results are proved for fields whose absolute Galois
group is a pro-$l$ group, although we expect to deal with the general case in
future work. Moreover, the evaluation of the Mackey functor completion would
be expected in general to be very difficult, since very simple questions
about, for example, derived functors of tensor products are quite difficult in
the Green functor situation [40]. C. Barwick has developed methods for proving
the Mackey/Green functor version of the conjecture. Presumably one will find
that our methods will permit the computation of the derived completions in the
Mackey/Green functor version, in terms of derived completions of ordinary
rings.
The paper is organized as follows. Section 2 records preliminary material,
including a discussion of actions of profinite groups on affine schemes which
will be important for the construction of ${\cal E}{G_{F}}$. Section 3 defines
totally torsion free profinite groups, and proves that absolute Galois groups
all enjoy this property. Section 4 develops a method for constructing
profinite group actions on certain (non-Noetherian) affine schemes. The
schemes in question are the inverse limits of systems of the form
$\cdots\stackrel{{\scriptstyle\times
l}}{{\longrightarrow}}\mathbb{G}_{m}^{n}\stackrel{{\scriptstyle\times
l}}{{\longrightarrow}}\mathbb{G}_{m}^{n}\stackrel{{\scriptstyle\times
l}}{{\longrightarrow}}\mathbb{G}_{m}^{n}\stackrel{{\scriptstyle\times
l}}{{\longrightarrow}}\mathbb{G}_{m}^{n}$
and ultimately ${\cal E}{G_{F}}$ will be constructed as an infinite product of
such schemes. We are able to construct actions on such varieties from affine
representations over $\mathbb{Z}_{l}$ of $G_{F}$. Section 5 actually
constructs ${\cal E}{G_{F}}$, using the totally torsion free property to
construct sufficiently many such affine representations. Section 6 is a
technical section, which studies the representation rings of profinite
$l$-groups which are an extension of a torsion free, topologically finitely
generated, abelian profinite $l$-group by a finite $l$-subgroup of the
symmetric group. Sections 7 and 8 now use this preparatory work to prove the
main result.
The author wishes to thank B. Conrad, P. Diaconis, D. Dugger, L. Hesselholt,
T. Lawson, G. Lyo, I. Madsen, D. Ramras, R. Vakil, and K. Wickelgren for
numerous helpful conversations on this work.
## 2 Preliminaries
### 2.1 General
We will assume that the reader is familiar with the theory of ring spectra and
modules over them, as well as the notion of a commutative ring spectrum. We
will use the version of this theory which uses symmetric spectra, which is
presented in [33]. In particular, we will assume the reader is familiar with
the standard choices of model structures on this category, and the associated
notions of cofibrant objects, as well as $Hom_{A}(M,N)$,
$M\mathrel{\mathop{\wedge}\limits_{A}}N$ and universal coefficient and Künneth
spectral sequences converging to them. We will also assume familiarity with
homotopy limits and colimits in this category. Throughout the paper, a
spectrum will mean a symmetric spectrum, a ring spectrum will mean a symmetric
ring spectrum, and a commutative ring spectrum will mean a commutative
symmetric ring spectrum.
### 2.2 $K$-theory
There are numerous versions of the construction of the $K$-theory spectra
associated with various kinds of categorical input data. We will require two
of them. The first is the construction to A. Elmendorf and M. Mandell [14]. It
takes as input small permutative categories, and produces from that
construction a spectrum. It advantage is that the category theoretic
information required to impose a commutative ring spectrum structure on the
resulting spectrum, and to produce module spectra over such spectra, is
codified in [14]. This kind of multiplicative structure is at the core of the
results of this paper. The second construction (due to Waldhausen [42]) takes
as input data a category with cofibrations and weak equivalences, and produces
a $K$-theory spectrum from this data. It has the advantage that it has the
excision properties proved to exist by Quillen ([31]), such as localization
and devissage. Each Waldhausen category can be regarded as a small permutative
category, to which the Elmendorff-Mandell construction applies, and there is a
sequence of natural equivalences of spectra relating the two constructions
which are described in [8], section 4. We will denote the Elmendorff-Mandell
construction by $\mathfrak{K}$ and the Waldhausen construction by $K$. We also
define a weakening of the notion of a ring spectrum as follows.
###### Definition 2.1
By a weak ring spectrum, we mean a spectrum $R$ equipped with a map $R\wedge
R\rightarrow R$ which is homotopy associative in the evident sense. By a
module over a weak ring spectrum $R$, we will mean a spectrum $M$ equipped
with a map $R\wedge M\rightarrow M$, where the diagrams involving the
structure maps commute up to homotopy. The Waldhausen construction readily
produces weak ring structures from tensor products on modules, as well as
modules from appropriate notions of tensor products.
Remark: The reason we need both constructions is the following. In order to
define our derived completion construction, we will need to have the precise
commutative ring structure on our spectra afforded by the Elmendorf-Mandell
model. On the other hand, for computation behavior of homotopy groups, it is
particularly important to realize constructions such as homotopy fibers as the
$K$-theory of categories themselves, and this can be carried out effectively
using the Waldhausen construction. We note that the structure of a weak ring
spectrum is sufficient to compute behavior on homotopy groups. The natural
equivalences allow us to go back and forth between models to obtain the
required information.
###### Proposition 2.1
Let $\underline{C}$ an abelian category, equipped with a strictly associative
symmetric monoidal product, which turns $\underline{C}$ into a bipermutative
category in the sense of [14]. Then the following hold.
1. 1.
$\mathfrak{K}(\underline{C})$ is naturally a commutative ring spectrum in the
sense of [14].
2. 2.
Let $\underline{D}\subseteq\underline{C}$ be a Serre subcategory, with the
property that it is closed under tensor product with any object in
$\underline{C}$. Then the quotient category $\underline{C}/\underline{D}$ is
also bipermutative, and $\mathfrak{K}(\underline{C}/\underline{D})$ is also a
commutative ring spectrum.
3. 3.
The natural equivalence between $\mathfrak{K}$ and $K$ respects the weak ring
spectrum structures, where the weak ring spectrum structure on $\mathfrak{K}$
is simply the forgetful one obtained from the full ring spectrum structure.
4. 4.
Given $\underline{C}$ and $\underline{D}$ as above, the homotopy fiber of the
evident map
$\mathfrak{K}(C)\rightarrow\mathfrak{K}(\underline{C}/\underline{D})$ is a
$\mathfrak{K}(\underline{C})$-module. As a weak $K(\underline{C})$-module, it
is equivalent to $K(\underline{D})$.
### 2.3 Equivariant algebraic $K$-theory
We will fix an algebraically closed ground field $k$ throughout this paper. In
this section we will define the equivariant algebraic-$K$-theory of finite
group actions on certain commutative rings. We will be addressing the fact
that in order to build the schemes we require, we will have to leave the
Noetherian setting. The suitable setting will be that of coherent commutative
rings, for which the category of finitely presented modules is an abelian
category. We will let $R[G]$ denote the representation ring of the group $G$
in an algebraically closed field $k$, which is a fixed ground field throughout
the paper. We also let $I_{G}$ denote the augmentation ideal, i.e. the kernel
of the natural augmentation $R[G]\rightarrow\mathbb{Z}$. We recall from [19]
the following definition.
###### Definition 2.2
A commutative ring $A$ is said to be coherent if for every homomorphism
$f:F_{0}\longrightarrow F_{1}$ of finitely generated free $A$-modules, the
kernel of $f$ is a finitely generated $A$-module.
A standard fact about coherent rings is the following.
###### Proposition 2.2
The category of finitely presented modules over a coherent ring $A$ is an
abelian category.
The following result is an extension of Proposition 4.5 of [8].
###### Proposition 2.3
Let $D$ be a directed set, and let $F$ be a functor from the associated
category to the category of rings. Suppose further that $F(d\leq d^{\prime})$
is faithfully flat for all $d\leq d^{\prime}$, and that all the rings $F(d)$
are coherent. Then the ring $\mathrel{\mathop{colim}\limits_{D}}F$ is
coherent.
Proof: The proof is identical to the proof of Proposition 4.5 of [8], and we
omit it. $\square$
###### Definition 2.3
We say a ring is Noetherian approximable if it can be described as a colimit
over a directed set of Noetherian rings, with the transition maps in the
colimit being faithfully flat.
Given a coherent commutative ring $A$, it is not known whether or not a
polynomial ring $A[x_{1},\ldots,x_{n}]$ is also coherent. However, this does
hold for Noetherian approximable rings $A$.
###### Proposition 2.4
Let $A$ be a Noetherian approximable commutative ring. Then the ring
$B=A[x_{s};s\in S]$, where $S$ is some indexing set, is Noetherian
approximable and therefore coherent.
Proof: $B$ can be described as the colimit of rings of the form
$R[x_{1},\ldots,x_{t}]$, where $R$ is Noetherian, and where the transition
maps in the colimit are faithfully flat. $R[x_{1},\ldots,x_{n}]$ is
Noetherian, and the maps in the system are clearly faithfully flat, hence the
result. $\square$
We let $A$ be a coherent commutative ring, equipped with an action by a finite
group $G$.
###### Definition 2.4
By a $G$-twisted $A$-module, we will mean a finitely presented $A$-module $M$
equipped with a $k$-linear $G$-action so that $g\cdot(\alpha\cdot
m)=\alpha^{g}\cdot(g\cdot m)$ for all $g\in G,\alpha\in A$, and $m\in M$. The
$G$-twisted $A$-modules form a category whose morphisms are the equivariant
$A$-linear maps, which we denote $\mbox{\em Mod}^{G}(A)$.
###### Proposition 2.5
Given $A$ and $G$ as above, the category $\mbox{\em Mod}^{G}(A)$ is an abelian
category.
Proof: Clear. $\square$
Although we will primarily be considering affine schemes, we will use ring and
scheme terminology interchangeably.
###### Definition 2.5
We will define the $G$-equivariant $K$-theory on the category of coherent
commutative rings to be the algebraic $K$-theory of the full subcategory of
$\mbox{Mod}^{G}(A)$ on the finitely presented $G$-twisted modules, and denote
it by $K^{G}(A)$. $K^{G}$ is functorial for flat homomorphisms of coherent
commutative rings with $G$ action. In the case of Noetherian rings, this is
what would usually be referred to as $G$-theory or $K^{\prime}$-theory.
In the case of a point, i.e. $A=k$, we have that $K^{G}(k)$ is simply the
category of finite dimensional $k$-linear representations of $G$.
$\pi_{0}(K^{G}(k))$ is then the $k$-linear representation ring $R[G]$. In
particular, if either (1) $Char(k)=0$ or (2) $Char(k)$ is relatively prime to
$\\#(G)$, then $\pi_{0}K^{G}(k)$ is isomorphic to the usual complex
representation ring or character ring of $G$. $K^{G}(k)$ is a commutative ring
spectrum. It is also equipped with an augmentation
$\epsilon:K^{G}(k)\rightarrow K(k)$, which forgets the $G$ action, and for
which applying $\pi_{0}$ induces the usual augmentation on representation
rings, which we also refer to as $\epsilon$.
We need a brief discussion of orbit schemes. From [21], it is known that for
finite group actions on affine schemes, there is always a categorical orbit
scheme which is the spectrum of the fixed point subring of the action.
###### Proposition 2.6
Suppose that the action of $G$ on an affine scheme $X=Spec(A)$ is free in the
sense that it is free on the set of all points, where a point denotes a
morphism $Spec(\Omega)\rightarrow X$ for some algebraically closed field
$\Omega$ containing $k$. Then there is a natural equivalence $K^{G}(A)\cong
K(A^{G})$. It is induced by the composite
$K(A^{G})\rightarrow K^{G}(A^{G})\rightarrow K^{G}(A)$
Equivalently, for free $G$-actions we have an equivalence $K(S/G)\rightarrow
K^{G}(X)$.
Proof: The equivalence of categories from which this equivalence results is
proved in [22]. $\square$
We call group actions satisfying the above conditions etale. The following is
well known.
###### Proposition 2.7
The category of affine schemes admits arbitrary inverse limits. The ring of
all global sections of an inverse limit is the corresponding colimit of the
rings of global sections of the component schemes.
###### Proposition 2.8
Let $G$ be a finite group, and let $\underline{X}$ denote a diagram of affine
schemes with $G$-action. If all the $G$ actions composing the diagram are
etale, then so is the $G$-action on $lim\underline{X}$.
Proof: For any algebraically closed field $\Omega$, the set of maps
$Spec(\Omega)\rightarrow lim\underline{X}$ is the inverse limit of the sets of
maps from $Spec(\Omega)$ to the component schemes. An inverse limit of free
$G$-sets over a directed set is always free. $\square$
Given an action of a finite group $G$ on a scheme $X$, and a $G$-invariant
subscheme $Y\subseteq X$, we define the category $\mbox{Nil}^{G}(X,Y)$ to be
the category of $G$-twisted $X$-modules whose support lies entirely in $Y$.
### 2.4 Profinite group actions on commutative rings
We begin with the definition of profinite groups. The basic material on
profinite groups is contained in [32].
###### Definition 2.6
A topological group $G$ is a profinite group if it is Hausdorff, compact, and
totally disconnected. It follows from this definition that $G$ is the inverse
limit of the inverse system of finite groups $\\{G/N\\}_{N}$, where the
indexing ranges over the closed normal subgroups of $G$ of finite index.
###### Proposition 2.9
Let $G$ be a profinite group. Then a subgroup of $G$ is open only if it is
closed and of finite index.
We next consider continuous actions of a profinite group $G$ on discrete sets
$X$.
###### Proposition 2.10
Let $X$ be a discrete set, and suppose that we are given an action of a
profinite group $G$ on $X$. Then the action is continuous if and only if the
stabilizers $G_{x}$ are closed subgroups of finite index of $G$ for all $x\in
X$.
Proof: From the discreteness of $X$, it follows that $\alpha^{-1}(x)$ is an
open set in $G\times X$ for all $x\in X$, where $\alpha:G\times X\rightarrow
X$ is the action map. A subset of $G\times X$ is open if and only if it is of
the form
$\coprod_{x\in X}U_{x}\times\\{x\\}$
where each $U_{x}$ is an open set in $G$. On the other hand,
$\alpha^{-1}(x)=\coprod_{\gamma\in\Gamma}G_{x}\gamma^{-1}\times\gamma\cdot x$
It is now clear that this set is open if and only if $G_{x}$ is an open
subgroup of $G$, i.e. that it is closed and of finite index by Proposition
2.9. $\square$
We also have the following result on transitive $G$-spaces.
###### Proposition 2.11
Let $X$ be a topological space with a continuous transitive action of a
profinite group $G$, and let $G_{x}$ denote the stabilizer of a point $x\in
X$. Then if $X$ is Hausdorff, we have that $X$ is $G$-homeomorphic to the
quotient space $G/G_{x}$, and $G_{x}$ is a closed subgroup of $G$.
Proof: It is immediate that $X$ is compact, since it is the surjective image
of the compact space $G$. The natural map $G/G_{x}\rightarrow X$ is a
bijective continuous map from one compact Hausdorff space to another, and is
therefore a homeomorphism. It is direct that $G/G_{x}$ is Hausdorff if and
only if $G_{x}$ is closed. $\square$
###### Definition 2.7
Let $G$ be a profinite group. By a continuous affine $G$-scheme $A$, we will
mean a continuous action of $G$ on a commutative ring $A$ equipped with the
discrete topology.
Example: Let $G$ and $G_{sep}$ be the Galois groups of $\overline{k}$, the
algebraic closure of a field $k$, and $k_{sep}$, its absolute closure. Then
the actions of $G$ and $G_{sep}$ on $Spec(\overline{k})$ and $Spec(k_{sep})$
are continuous affine $G$ (respectively $G_{sep}$)-schemes.
Example: Let $G$ denote the $l$-adic integers, and let
$A=\bigcup_{n}k[t^{\pm\frac{1}{l^{n}}}]$. The group of automorphisms of $A$
leaving $k[t^{\pm 1}]$ elementwise fixed is isomorphic to $\mathbb{Z}_{l}$,
and the action of $\mathbb{Z}_{l}$ on $A$ is a continuous affine
$\mathbb{Z}_{l}$-scheme.
Let $G$ be a profinite group, and let $A$ be a continuous affine $G$-scheme.
Let $B$ denote the fixed point subring $A^{G}$. For every closed normal
subgroup $N$ of finite index in $G$, we may consider the scheme $Spec(A^{N})$.
These schemes fit together into an inverse system of affine schemes
$\\{Spec(A^{N})\\}_{N}$, and we may regard the inverse system as an inverse
system of sets, by simply considering the points of the schemes. We also have
the group actions
$G/N\times Spec(A^{N})\longrightarrow Spec(A^{N})$
which fit together into an action of a pro-group on a pro-set. For every
$\mathfrak{p}\in Spec(B)$, we may consider the fiber over $\mathfrak{p}$ in
$\mathrel{\mathop{lim}\limits_{\longleftarrow}}Spec(A^{N})$, denoted by
$F_{\mathfrak{p}}$, and find that we obtain a continuous action of $G$ on
$F_{\mathfrak{p}}$.
###### Proposition 2.12
Let $A$ be a continuous affine $G$-scheme, with $B=A^{G}$. Then the map
$\pi:Spec(A)\rightarrow Spec(B)$
is surjective, and all the fibers $\pi^{-1}(\mathfrak{p})$ are transitive
$G$-sets. Moreover, $\pi^{-1}(\mathfrak{p})$ is homeomorphic to a $G$-space of
the form $G/G^{\prime}$, where $G^{\prime}$ is a closed subgroup of $G$.
Proof: For every closed subgroup of finite index $K\subseteq G$, we have the
subring $A^{K}$ of elements of $A$ fixed by $K$. The ring $A$ is equal to the
union $\bigcup A^{K}$ as $K$ varies over smaller and smaller closed subgroups
of finite index in $G$, because the stabilizer of every element of $A$ is a
closed subgroup of finite index in $G$ by Proposition 2.10. We will now show
that the natural map
$\alpha:Spec(A)\longrightarrow\mbox{
}\mathrel{\mathop{lim}\limits_{\longleftarrow}}Spec(A^{K})$
is a bijection. To show that $\alpha$ is injective, consider two prime ideals
$\mathfrak{p}$ and $\mathfrak{q}$ in $A$. If $\mathfrak{p}\neq\mathfrak{q}$,
then there is an element $a\in A$ with $a$ an element of one of $\mathfrak{p}$
and $\mathfrak{q}$ and not in the other, say $a\in\mathfrak{p}$ and
$a\notin\mathfrak{q}$. But we have seen that $a\in A^{K}$ for some closed
subgroup $K$ of finite index in $G$. It follows that the images of
$\mathfrak{p}$ and $\mathfrak{q}$ in $Spec(A^{K})$ are distinct, which gives
the injectivity of $\alpha$. For surjectivity, we first note that elements
$\mathrel{\mathop{lim}\limits_{\longleftarrow}}Spec(A^{K})$ correspond exactly
to families $\\{\mathfrak{p}_{K}\\}_{K}$, where the system varies over the
closed subgroups of finite index, where each $\mathfrak{p}_{K}$ is a prime
ideal in $A^{K}$, and where we have the consistency condition
$\pi^{K}_{K^{\prime}}(\mathfrak{p}_{K})=\mathfrak{p}_{K^{\prime}}$ whenever
$K\subseteq K^{\prime}$. Here $\pi^{K}_{K^{\prime}}$ denotes the restriction
map $Spec(A^{K})\rightarrow Spec(A^{K^{\prime}})$. Given a family ${\cal
F}=\\{\mathfrak{p}_{k}\\}_{K}$ as above, we define a subset
$\mathfrak{p}_{{\cal F}}\subseteq A$ to be the set of all elements $a\in A$
such that there exists a closed subgroup $K\subseteq G$ of finite index so
that $a\in\mathfrak{p}_{K}\subseteq A^{K}$. We will verify that
$\mathfrak{p}_{{\cal F}}$ is a prime ideal in $A$. We first make the
elementary observation that if $a\in\mathfrak{p}_{K}$ for some $K$, then
$a\in\mathfrak{p}_{K^{\prime}}$ for all closed subgroups of finite index
$K^{\prime}\subseteq K$. We verify that $\mathfrak{p}_{{\cal F}}$ is closed
under addition. Suppose $a,a^{\prime}\in\mathfrak{p}_{{\cal F}}$, and fix
closed subgroups of finite index $K$ and $K^{\prime}$ so that
$a\in\mathfrak{p}_{K}$ and $a^{\prime}\in\mathfrak{p}_{K^{\prime}}$. Now, the
subgroup $K\cap K^{\prime}$ is also a closed subgroup of finite index, and
both $a$ and $a^{\prime}$ are elements of $\mathfrak{p}_{K\cal K^{\prime}}$,
consequently $a+a^{\prime}\in\mathfrak{p}_{K\cal K^{\prime}}$, which means
that $a+a^{\prime}\in\mathfrak{p}_{{\cal F}}$. To see that
$\mathfrak{p}_{{\cal F}}$ is an ideal, we consider $a\in A$ and
$i\in\mathfrak{p}_{{\cal F}}$. Select closed subgroups of finite index $K$ and
$K^{\prime}$ so that $a\in A^{K}$ and $i\in\mathfrak{p}_{K^{\prime}}$. We then
have $a\in A^{K\cap K^{\prime}}$ and $i\in\mathfrak{p}_{K\cap K^{\prime}}$,
and therefore that $a\cdot i\in\mathfrak{p}_{K\cap K^{\prime}}$. Finally, to
see that $\mathfrak{p}_{{\cal F}}$ is prime, consider $a,a^{\prime}\in A$ for
which $a\cdot a^{\prime}\in\mathfrak{p}_{{\cal F}}$. Select a closed subgroup
of finite index $K$ so that $a,a^{\prime}\in A^{K}$. Since $a\cdot
a^{\prime}\in\mathfrak{p}_{{\cal F}}$, there is another closed subgroup of
finite index $K^{\prime}$ so that $a\cdot
a^{\prime}\in\mathfrak{p}_{K^{\prime}}$. We then have that $a,a^{\prime}$, and
$a\cdot a^{\prime}$ all lie in $A^{K\cap K^{\prime}}$ and that $a\cdot
a^{\prime}\in\mathfrak{p}_{K\cap K^{\prime}}$. Because of the primeness of
$\mathfrak{p}_{K\cap K^{\prime}}$, it follows that either $a$ or $a^{\prime}$
is an element of $\mathfrak{p}_{K\cap K^{\prime}}$, and consequently that
either $a$ or $a^{\prime}$ lies in $\mathfrak{p}_{{\cal F}}$. Under the
identification
$Spec(A)\longrightarrow\mathrel{\mathop{lim}\limits_{\longleftarrow}}Spec(A^{N})$
The fiber $\pi^{-1}(\mathfrak{p})$ is identified with $F_{\mathfrak{p}}$
above.
We next note that the maps $\pi^{K}_{K^{\prime}}:Spec(A^{K})\rightarrow
Spec(B)$ are surjective for any closed normal subgroup of finite index
$K\subseteq G$. This result is immediate from Proposition 1.1 of [21],
applying it to the action of the finite group $G/K$ on $Spec(A^{K})$. The
surjectivity statement above now follows immediately.
For the transitivity of the $G$ action on the fibers of the map $\pi$, we
proceed as follows. We first note that $Spec(A)$ can also be described as
$\mathrel{\mathop{lim}\limits_{\longleftarrow}}Spec(A^{N})$
where $N$ ranges over the closed normal subgroups of finite index in $G$, due
to the finality of the closed normal subgroups of finite index in the set of
all closed subgroups of finite index. For each such $N$, and pair of elements
$x,y\in Spec(A)$, we consider the subset $\Phi_{N}(x,y)\subseteq G$ consisting
of all $g\in G$ such that $gx_{N}=y_{N}$, where $x_{N}$ and $y_{N}$ denote the
projections of $x$ and $y$ in $Spec(A^{N})$. We are letting $G$ act on
$Spec(A^{N})$ in the obvious way through the projection $G\rightarrow G/N$.
The set $\Phi_{N}(x,y)$ is non-empty due to the transitivity statement for
finite group actions on rings made in Proposition 1.1 of [21]. Moreover, if we
have $N^{\prime}\subseteq N$, then
$\Phi_{N^{\prime}}(x,y)\subseteq\Phi_{N}(x,y)$. It is clear that the sets
$\Phi_{N}(x,y)$ are closed subsets of $G$, and for a fixed choice of $x$ and
$y$, have the finite intersection property. Therefore, due to the compactness
of $G$,
$\bigcap_{N}\Phi_{N}(x,y)\neq\emptyset$
Any element in $\bigcap_{N}\Phi_{N}(x,y)$ will now carry $x$ to $y$, giving
the transitivity statement. We now know that $\pi^{-1}(\mathfrak{p})$ is a
Hausdorff transitive continuous $G$-set, and conclude from Proposition 2.11
that $\pi^{-1}(\mathfrak{p})$ is $G$-homeomorphic to a $G$-space of the form
$G/G^{\prime}$, where $G^{\prime}$ is a closed subgroup of $G$. $\square$
###### Corollary 2.1
For continuous affine $G$-schemes, a categorical orbit space exists, and
$Spec(A^{G})$ is a choice of such an orbit space.
Proof: This is proved as in the finite case, performed in [21]. $\square$
Let $A$,$B$, and $G$ be as above. Let $\mathfrak{p}$ be a prime ideal of $A$,
and let $G_{\mathfrak{p}}$ denote its stabilizer, which will be called its
decomposition group. $G_{\mathfrak{p}}$ is a closed subgroup of $G$. Let
$k({\mathfrak{p}})$ denote the field of fractions of $A/\mathfrak{p}$.
Similarly, let $\mathfrak{q}=\mathfrak{p}\cap B$, and write $k(\mathfrak{q})$
for the field of fractions of $B/\mathfrak{q}$. There is an evident inclusion
$k(\mathfrak{q})\subseteq k(\mathfrak{p})$. The group $G_{\mathfrak{p}}$ acts
continuously on $k(\mathfrak{p})$, when $k(\mathfrak{p})$ is equipped with the
discrete topology. There is therefore a homomorphism from $G_{\mathfrak{p}}$
to the Galois group $Gal(k(\mathfrak{p}),k(\mathfrak{q}))$, which is in
general a profinite group. This homomorphism is surjective. This is stated
explicitly in [21] in the case where $G$ is finite, and the result in this
context follows immediately from this case by passing to inverse limits over
Hausdorff finite quotients. The kernel of this homomorphism is called the
inertia group of $\mathfrak{p}$, and is clearly a closed subgroup of
$G_{\mathfrak{p}}$. We denote it by $I_{\mathfrak{p}}$.
We will also study the notion of etale actions in the context of profinite
actions.
###### Definition 2.8
Let $G$ be a profinite group, and let $A$ be a continuous affine $G$-scheme.
We say the action is etale if each of the inclusions $A^{G}\hookrightarrow
A^{K}$ is an etale map for every closed subgroup of finite index $K$. This is
clearly equivalent to the condition that each inclusion $A^{G}\hookrightarrow
A^{N}$ is a Galois cover with group $G/N$, where $N$ varies over all the
closed normal subgroups of finite index.
We have the following analogue of Corollary 2.4 of [21] for verifying that an
action is etale.
###### Proposition 2.13
Let $A$ be a commutative ring equipped with the discrete topology, and let a
profinite group $G$ act continuously on $A$. Then the action is etale if and
only if the inertia group $I_{\mathfrak{p}}$ is trivial for all prime ideals
$\mathfrak{p}$ of $A$.
Proof: This is a straightforward consequence of the case of a finite group
$G$, for which a proof is given in §2 of [21]. $\square$
There is a version of this criterion which will be useful to us. Let $A$ and
$G$ be as above, and let $\mathfrak{q}$ denote any prime ideal of $A^{G}$. Let
$\Omega$ denote any algebraically closed field which admits an embedding
$i:k(\mathfrak{q})\rightarrow\Omega$. We let $\mathfrak{H}(A,\Omega,i)$ denote
the set of all ring homomorphisms $A\rightarrow\Omega$ extending the composite
$A^{G}\rightarrow A^{G}/\mathfrak{q}\hookrightarrow
k(\mathfrak{q})\stackrel{{\scriptstyle i}}{{\longrightarrow}}\Omega$
For any closed normal subgroup $N$ of $G$ of finite index, we similarly have
the set $\mathfrak{H}(A^{N},\Omega,i)$ and we have restriction maps
$\mathfrak{H}(A^{N^{\prime}},\Omega,i)\longrightarrow\mathfrak{H}(A^{N},\Omega,i)$
whenever $N^{\prime}\subseteq N$. It is readily verified that
$\mathfrak{H}(A,\Omega,i)\cong\mathrel{\mathop{lim}\limits_{\leftarrow}}\mathfrak{H}(A^{N},\Omega,i)$
and therefore further that $\mathfrak{H}(A,\Omega,i)$ is equipped with a
continuous $G$ action.
###### Proposition 2.14
The $G$-space $\mathfrak{H}(A,\Omega,i)$ is transitive, and is
$G$-homeomorphic to the quotient $G/I_{\mathfrak{p}}$, where $\mathfrak{p}$ is
the kernel of some ring homomorphism $A\rightarrow\Omega$ extending $i$.
Proof: For the finite case, this is proved in §2 of [21]. The extension to the
profinite case is direct. $\square$
###### Corollary 2.2
A continuous affine $G$-scheme is etale if and only if the $G$ actions on all
sets $\mathfrak{H}(A,\Omega,i)$ are free, for all prime ideals $\mathfrak{p}$
of $A^{G}$ and homomorphisms $i:k(\mathfrak{p})\hookrightarrow\Omega$ for all
algebraically closed fields $\Omega$.
In order to define the equivariant algebraic $K$-theory of continuous affine
$G$-schemes, we extend the quasi-Noetherian condition to such actions.
###### Definition 2.9
We say that a continuous affine $G$-scheme $(A,G)$ is coherent if $A$ is a
coherent ring. We say it is totally coherent if $A^{N}$ is a coherent ring for
all closed subgroups $N\subseteq G$.
We now define the category of twisted $G$ modules for coherent continuous
affine $G$-schemes.
###### Definition 2.10
Let $(A,G)$ be a ring $A$ equipped with a continuous action by a profinite
group $G$, where we assume that $A$ is equipped with the discrete topology. By
a $G$-twisted module over $A$, we mean an $A$ -module $M$, endowed with the
discrete topology, together with a continuous action by $G$, which is
compatible with the $G$ action on $A$ via the formula
$g\cdot(\alpha m)=\alpha^{g}\cdot gm$
The $G$-twisted $A$-modules form an abelian category. A $G$-twisted $A$-module
is said to be finitely presented if the underlying module is a finitely
presented $A$-module. In the case of coherent continuous affine $G$-schemes,
the finitely presented $G$-twisted $A$-modules also form an abelian category.
We will write $\mbox{\em Mod}^{G}(A)$ for the category of finitely presented
$G$-twisted $A$-modules.
We now define the equivariant algebraic $K$-theory of coherent continuous
affine $G$-schemes.
###### Definition 2.11
Given a coherent continuous affine $G$-scheme $(A,G)$, by the equivariant
$K$-theory of $A$, denoted $K^{G}(A)$, we will mean the $K$-theory of the
Waldhausen category associated to the abelian category $\mbox{\em
Mod}^{G}(A)$.
With this definition, we would now like to provide an analogue to Proposition
2.6 for actions of a profinite group.
Begin with some results on coherent rings.
###### Proposition 2.15
Let $f:A\rightarrow B$ be a faithfully flat morphism of rings. Then an
$A$-module $M$ is finitely generated (respectively finitely presented) if and
only if $B\mathrel{\mathop{\otimes}\limits_{A}}M$ is finitely generated
(finitely presented).
Proof: This is Corollary 1.11 of [22] $\square$
###### Corollary 2.3
Let $f:A\rightarrow B$ be a faithfully flat morphism of rings. If $B$ is
coherent, then so is $A$.
Proof: We will need to prove that the kernel of a homomorphism $f:F\rightarrow
F^{\prime}$ between finitely generated $A$-modules is finitely generated.
Because of the flatness, we have an exact sequence
$0\rightarrow B\mathrel{\mathop{\otimes}\limits_{A}}Ker(f)\rightarrow
B\mathrel{\mathop{\otimes}\limits_{A}}F\rightarrow
B\mathrel{\mathop{\otimes}\limits_{A}}F^{\prime}$
The coherence of $B$ now gives that
$B\mathrel{\mathop{\otimes}\limits_{A}}Ker(f)$ is finitely generated.
Proposition 2.15 now gives the result. $\square$
###### Corollary 2.4
Let $A$ be a coherent commutative ring equipped with the discrete topology,
and let a profinite group $G$ act continuously on $A$. If the $G$-action is
etale, then the fixed point subring $A^{G}$ is coherent.
Proof: The action being etale implies that the inclusion $A^{G}\hookrightarrow
A$ is faithfully flat. $\square$
###### Corollary 2.5
For any continuous etale action of a profinite group $G$ on a coherent ring
$A$, $A^{G}$ is coherent there is a natural map
$K(A^{G})\rightarrow K^{G}(A^{G})\rightarrow K^{G}(A)$
In the terminology of schemes $X=Spec(A)$, there is a map
$\eta_{X}:K(X/G)\rightarrow K^{G}(X)$ when $X$ is a coherent affine scheme,
natural for flat maps of schemes.
We wish to prove that $\eta_{X}$ is an equivalence. In the case of a finite
group $G$, this result follows from Theorem II.5.1 of [25], by proving that
the functors $M\rightarrow M^{G}$ from $\mbox{Mod}^{G}(A)$ to
$\mbox{Mod}^{\\{e\\}}(A^{G})$ and $M\rightarrow
B\mathrel{\mathop{\otimes}\limits_{A}}M$ from $\mbox{Mod}^{\\{e\\}}(A^{G})$ to
$\mbox{Mod}^{G}(A)$ are inverse equivalences of categories. The fact that both
functors preserve the subcategories of finitely presented modules follows from
the faithful flatness of the inclusion $A^{G}\hookrightarrow A$ together with
Proposition 2.15 above.
###### Proposition 2.16
The natural transformation $\eta_{X}$ is an equivalence for any continuous
etale action of a profinite group $G$ on an affine scheme $X=Spec(A)$, where
$A$ is coherent.
Proof: We may construct the categories $\mbox{Mod}^{G/N}(A^{N})$ for any
closed normal subgroup $N\subseteq G$ of finite index. We note that all the
rings $A^{N}$ are coherent by Corollary 2.3, and that all $N\subseteq
N^{\prime}$ the inclusions $A^{N}\hookrightarrow A^{N^{\prime}}$ are
faithfully flat, because they are etale ring extensions. It follows that we
have a directed system $N\rightarrow\mbox{Mod}^{G/N}(A^{N})$, where the
functors
$\mbox{Mod}^{G/N}(A^{N})\rightarrow\mbox{Mod}^{G/N^{\prime}}(A^{N^{\prime}})$
are defined to be $A^{N^{\prime}}\mathrel{\mathop{\otimes}\limits_{A^{N}}}-$.
There is therefore a functor
$j:\mathrel{\mathop{colim}\limits_{\mathrel{\mathop{\longrightarrow}\limits_{N}}}}\mbox{Mod}^{G/N}(A^{N})\rightarrow\mbox{Mod}^{G}(A)$
which we claim is an equivalence of categories. This will suffice to prove the
result, because each of the functors
$\mbox{Mod}^{\\{e\\}}(A^{G})\rightarrow\mbox{Mod}^{G/N}(A^{N})$ is an
equivalence of categories by the validity of the result for finite group
actions mentioned above.
We proceed to prove that $j$ is an equivalence of categories. We first observe
that $j$ is surjective on isomorphism classes of objects. In order to do this,
we will first give a description of $G$-twisted modules over $A$ whose
underlying $A$-module is free in terms of matrices. Suppose $F$ is a free
module with continuous $G$-action twisted compatibly with the action of $G$ on
$A$. We fix a basis $B=\\{b_{1},\ldots,b_{t}\\}$ for $F$. Because of the
continuity of the action, it is clear that there is a closed normal subgroup
$N$ of finite index in $G$ so that $N$ fixes the set $B$ elementwise. It
follows that the action of an element $g\in G$ on an element of $B$ depends
only on its coset in $G/N$. For each coset $\gamma$ in $G/N$, we then obtain a
$t\times t$ matrix $M(\gamma)$ with entries in $A$, defined by
$M(\gamma)(b_{j})=\sum_{i}M(\gamma)_{ij}b_{i}$. Let $N$ be such that
$M(\gamma)_{ij}\in A^{N}$ for all $\gamma$ and pairs $(i,j)$. The coefficients
of all the matrices $M(\gamma)$ are therefore acted on trivially by $N$, and
therefore the action of $G$ on the free $A^{N}$-submodule generated by $B$
factors through $G/N$. The matrices $M(\gamma)$ clearly satisfy the relations
$M(\gamma_{1}\gamma_{2})=M(\gamma_{1})\cdot M(\gamma_{2})^{\gamma_{1}}$ (2–1)
where the superscript indicates entrywise action of $G$ on the matrix
$M(\gamma_{2})$. Sets of matrices $\\{M_{\gamma}\\}_{\gamma}\in G$ satisfying
these relations in turn determine a $G/N$-twisted module over $A^{N}$, and
after applying $A\mathrel{\mathop{\otimes}\limits_{A^{N}}}-$, over $A$. It is
furthermore clear that the result is isomorphic to the original $G$-twisted
$A$-module, so we have proved that any continuous twisted $G$-module over $A$,
with free underlying module, is obtained by extension of scalars from a free
twisted $G/N$-module over $A$, for some closed normal subgroup of finite index
$N$.
Next, we will prove that given any homomorphism of $G$-twisted $A$-modules
$f:A\mathrel{\mathop{\otimes}\limits_{A^{N}}}F_{0}\rightarrow
A\mathrel{\mathop{\otimes}\limits_{A^{N}}}F_{1}$
where $F_{0}$ and $F_{1}$ are $G/N$-twisted $A^{N}$-modules whose underlying
modules are finitely generated and free, there is an $N^{\prime}\supseteq N$
and homomorphism of $G/N^{\prime}$-twisted $A^{N^{\prime}}$-modules
$\overline{f}:A^{N^{\prime}}\mathrel{\mathop{\otimes}\limits_{A^{N}}}F_{0}\rightarrow
A^{N^{\prime}}\mathrel{\mathop{\otimes}\limits_{A^{N}}}F_{1}$
so that $f=A\mathrel{\mathop{\otimes}\limits_{A^{N^{\prime}}}}\overline{f}$.
To see this, we observe that given bases for $F_{0}$ and $F_{1}$, $f$
determines a matrix $M_{f}$ with entries in $A$. The $G$-twisted modules
$F_{0}$ and $F_{1}$ determine families of matrices
$\\{M_{0}(\gamma)\\}_{\gamma\in G}$ and $\\{M_{1}(\gamma)\\}_{\gamma\in G}$
with entries in $A^{N}$, and the fact that $f$ is a homomorphism of
$G$-twisted $A$-modules means that the equations
$M_{1}({\gamma})M_{f}^{\gamma}=M_{f}M_{0}(\gamma)$ (2–2)
hold for all $\gamma\in G$, when $M_{0}(\gamma)$ and $M_{1}(\gamma)$ are
regarded as matrices with entries in $A$ via the homomorphism
$A^{N}\rightarrow A$. As above, there is a $N^{\prime}\supseteq N$ so that all
the entries of $M_{f}$ lie in the image of the homomorphism
$A^{N^{\prime}}\rightarrow A$, and we let $\overline{M}_{f}$ denote the matrix
$M_{f}$ regarded as a matrix with entries in $A^{N^{\prime}}$. We know that
the relations (2–2) hold in $A$, but this means that the relations
$M_{1}({\gamma})\overline{M}_{f}^{\gamma}=\overline{M}_{f}M_{0}(\gamma)$
hold in $A^{N^{\prime}}$ by the faithful flatness of the homomorphism
$A^{N^{\prime}}\rightarrow A$. The required homomorphism $\overline{f}$ is now
the homomorphism represented by the matrix $\overline{M}_{f}$.
Next consider any finitely presented $G$-twisted $A$-module $M$, with a
presentation
$F_{1}\stackrel{{\scriptstyle f}}{{\rightarrow}}F_{0}\rightarrow M\rightarrow
0$ (2–3)
where $F_{0}$ and $F_{1}$ are $G$-twisted $A_{F}$-modules with finitely
generated free underlying modules. Our first argument above shows that there
is closed normal subgroup of finite index $N$ so that $F_{0}$ and $F_{1}$ are
extended from $A^{N}$, and the second shows that there is a
$N^{\prime}\supseteq N$ so that $f$ is extended from $A^{N^{\prime}}$. Let
$\overline{F}_{0},\overline{F}_{1}$ be $G/N^{\prime}$-twisted
$A^{N^{\prime}}$-modules (with free underlying modules) and $\overline{f}$ be
a homomorphism $\overline{f}:\overline{F}_{1}\rightarrow\overline{F}_{0}$ of
$G/N^{\prime}$-twisted $A^{N^{\prime}}$-modules so that the presentation 2–3
is obtained by applying $A\mathrel{\mathop{\otimes}\limits_{A^{N^{\prime}}}}-$
to the exact sequence
$\overline{F}_{1}\stackrel{{\scriptstyle\overline{f}}}{{\rightarrow}}\overline{F}_{0}\rightarrow
C\rightarrow 0$
of $G/N^{\prime}$-twisted $A^{N^{\prime}}$-modules, where $C$ denotes the
cokernel of $\overline{f}$. It follows by flatness that $M\cong
A\mathrel{\mathop{\otimes}\limits_{A^{N^{\prime}}}}C$. It follows that the
functor $j$ above is surjective on isomorphism classes of objects.
To prove injectivity on isomorphism classes of objects, we suppose that we are
given a closed normal subgroup of finite index $N\subseteq G$ and finitely
presented $G/N$-twisted $A^{N}$-modules $M$ and $M^{\prime}$, so that
$A\mathrel{\mathop{\otimes}\limits_{A^{N}}}M$ and
$A\mathrel{\mathop{\otimes}\limits_{A^{N}}}M^{\prime}$ are isomorphic. We
suppose we are given an isomorphism
$f:A\mathrel{\mathop{\otimes}\limits_{A^{N}}}M\rightarrow
A\mathrel{\mathop{\otimes}\limits_{A^{N}}}M^{\prime}$ and its inverse
$f^{-1}$. Both $f$ and $f^{-1}$ are of the form
$A\mathrel{\mathop{\otimes}\limits_{A^{N^{\prime}}}}\overline{f}$ and
$A^{N^{\prime}}\mathrel{\mathop{\otimes}\limits_{A^{N}}}\overline{f^{-1}}$ for
some $N^{\prime}$ and homomorphisms
$\overline{f}:A^{N^{\prime}}\mathrel{\mathop{\otimes}\limits_{A^{N}}}M\rightarrow
A^{N^{\prime}}\mathrel{\mathop{\otimes}\limits_{A^{N}}}M^{\prime}$
and
$\overline{f^{-1}}:A^{N^{\prime}}\mathrel{\mathop{\otimes}\limits_{A^{N}}}M^{\prime}\rightarrow
A^{N^{\prime}}\mathrel{\mathop{\otimes}\limits_{A^{N}}}M$
But $\overline{f}$ and $\overline{f^{-1}}$ are also inverses, because they
become so over $A$, and because of the faithful flatness of
$A^{N^{\prime}}\rightarrow A$. It follows that $\overline{f}$ is an
isomorphism, which gives the injectivity result for isomorphism classes of
objects.
Finally we will need to check that for any pair of objects $x,y$ in
$\underline{C}=\mathrel{\mathop{colim}\limits_{\mathrel{\mathop{\longrightarrow}\limits_{N}}}}\mbox{Mod}^{G/N}(A^{N})$,
$j$ induces a bijection
$Hom_{\underline{C}}(x,y)\stackrel{{\scriptstyle\sim}}{{\rightarrow}}Hom_{\mbox{Mod}^{G}(A)}(j(x),j(y))$
The injectivity of this map is immediate from the faithful flatness of all the
maps $A^{N}\rightarrow A$. To see the surjectivity, we assume that the objects
$x$ and $y$ are represented by $G/N$-twisted $A^{N}$-modules $M$ and
$M^{\prime}$, respectively, equipped with presentations
$F_{1}\stackrel{{\scriptstyle\alpha}}{{\rightarrow}}F_{0}\rightarrow
M\rightarrow 0$
and
$F_{1}^{\prime}\stackrel{{\scriptstyle\beta}}{{\rightarrow}}F_{0}^{\prime}\rightarrow
M^{\prime}\rightarrow 0$
A homomorphism $f$ from $M$ to $M^{\prime}$ can be represented by a pair of
homomorphisms $f_{0}:F_{0}\rightarrow F_{0}^{\prime}$ and
$f_{0}:F_{1}\rightarrow F_{1}^{\prime}$ so that $f_{0}\alpha=\beta f_{1}$.
These homomorphisms can be represented by matrices, and by arguments similar
to the ones above we find that these matrices can be represented over $A^{N}$
for some closed normal subgroup of finite index $N\subseteq G$. This gives the
result. $\square$
### 2.5 Homotopy properties
Let $G$ be any finite group, and let $G$ act by automorphisms on a vector
space $V$ over a field $k$ via a representation $\rho$f. We suppose that
$\\#(G)$ is relatively prime to $char(k)$. We may form the skew polynomial
ring $k[V]_{\rho}[G]$ (see [20]). The following is now a formal observation.
###### Proposition 2.17
The category $\mbox{Mod}^{G}(V)$ ($V$ is regarded as an affine scheme) is
equivalent to the category of finitely generated modules over the regular
Noetherian ring $k[V]_{\rho}[G]$.
Let $\mbox{Rep}[G]$ denote the category of finite dimensional $k$-linear
representations of $G$, which can be identified with the category of finitely
generated modules over the group ring $k[G]$. We have an inclusion
$i:k[G]\hookrightarrow k[V]_{\rho}[G]$ and a projection
$\pi:k[V]_{\rho}[G]\rightarrow k[G]$ of rings, sending the vector space $V$
identically to $0$. The following result is now an immediate consequence of
Theorem 7 of [31].
###### Theorem 2.1
The ring homomorphisms $i$ and $\pi$ give rise to exact functors
$k[V]_{\rho}[G]\mathrel{\mathop{\otimes}\limits_{k[G]}}-:\mbox{\em
Rep}[G]\rightarrow\mbox{ \em Mod}^{G}(V)$
and
$k[G]\mathrel{\mathop{\otimes}\limits_{k[V]_{\rho}[G]}}-:\mbox{\em
Mod}^{G}(V)\rightarrow\mbox{\em Rep}[G]$
Both functors induce equivalences on $K$-theory spectra.
We will also need to describe how induction of representations works in this
context. Let $H\subseteq G$ be an inclusion of finite groups, and let $\rho$
denote a $k$-linear representation of $H$, with $V_{\rho}$ as the
representation space. The induced representation $i_{H}^{G}(\rho)$ can be
defined via the operation
$k[G]\mathrel{\mathop{\otimes}\limits_{k[H]}}-:\mbox{Mod}(k[H])\rightarrow\mbox{Mod}(k[G])$
where for a ring $A$ the category of finitely generated left $A$ -modules is
denoted by $\mbox{Mod}(A)$. We denote the representation space of
$i_{H}^{G}(\rho)$ by $V_{i_{H}^{G}(\rho)}$. There is a natural $H$-equivariant
inclusion $V_{\rho}\hookrightarrow V_{i_{H}^{G}(\rho)}$, and consequently an
inclusion of rings
$k[V_{\rho}]_{\rho}[H]\hookrightarrow
k[V_{i_{H}^{G}(\rho)}]_{{i_{H}^{G}(\rho)}}[G]$
###### Proposition 2.18
Given $H,G$, and $\rho$ as above, the following diagram of functors commutes.
$\begin{CD}\mbox{{\em
Mod}}(k[H])@>{\phantom{\mbox{xxxxxxx}}{k[G]\mathrel{\mathop{{\otimes}}\limits_{k[H]}}-}\phantom{\mbox{xxxxxxxxxx}}}>{}>\mbox{\em
Mod}(k[G])\\\
@V{k[V_{\rho}]_{\rho}[H]\mathrel{\mathop{\otimes}\limits_{k[H]}}-}V{}V@V{}V{k[V_{i^{G}_{H}(\rho)}]_{i_{H}^{G}(\rho)}[G]\mathrel{\mathop{\otimes}\limits_{k[G]}}-}V\\\
\mbox{\em
Mod}(k[V_{\rho}]_{\rho}[H])@>{k[V_{i^{G}_{H}(\rho)}]_{i_{H}^{G}(\rho)}[G]\mathrel{\mathop{\otimes}\limits_{k[V_{\rho}]_{\rho}[H]}}-}>{}>\mbox{\em
Mod}(k[V_{i^{G}_{H}(\rho)}]_{i_{H}^{G}(\rho)}[G])\end{CD}$
### 2.6 Completions
We recall from [7] that given a homomorphism of commutative ring spectra
$f:A\rightarrow B$, and an $A$-module $M$, one may construct the derived
completion $M^{\wedge}_{f}$. Here are some examples of this construction.
Example: Any commutative ring may be thought of as a commutative ring spectrum
via the Eilenberg-MacLane construction. If $f:A\rightarrow B$ is a surjective
homomorphism of commutative Noetherian rings, and $M$ is a finitely generated
$A$-module, then the derived completion $M^{\wedge}_{f}$ is the Eilenberg-
MacLane construction on the usual completion $M^{\wedge}_{I}$, where $I$ is
the kernel of $f$.
Example Let $f:S^{0}\rightarrow\mathbb{H}_{l}$ be the mod $l$ Hurewicz map,
and let $X$ be any spectrum (and therefore an $S^{0}$-module). Then
$X^{\wedge}_{f}$ is the usual Bousfield-Kan notion of completion of the
spectrum at the prime $l$ (see [6]).
Example: Let $A=R[\mathbb{Z}_{l}]$ be the representation ring of the additive
group of the $l$-adic integers, defined as the direct limit of the
representation rings of the finite cyclic $l$-th power order groups. Let
$f:A\rightarrow\mathbb{F}_{l}$ be the augmentation followed by reduction mod
$l$. Then $A^{\wedge}_{f}$ is equivalent to the $l$-adic completion of the
integral group ring of the simplicial group of points on the circle group. In
particular, the homotopy groups are $\cong\mathbb{Z}_{l}$ for dimensions 0 and
1, and are $\cong 0$ otherwise.
It has a number of useful properties which can be found in [7], and we record
some of them here. Let $f:A\rightarrow B$ be a homomorphism of commutative
ring spectra.
###### Proposition 2.19
Suppose that
$M\rightarrow N\rightarrow P$
is a cofibration sequence of $A$-modules. Then the naturally defined sequence
$M^{\wedge}_{f}\rightarrow N^{\wedge}_{f}\rightarrow P^{\wedge}_{f}$
is also a cofibration sequence.
###### Proposition 2.20
Suppose we are given a homomorphism $\lambda:M\rightarrow N$ of $A$-modules.
Suppose further that the derived smash product map
$id_{B}\wedge\lambda:B\mathrel{\mathop{\wedge}\limits_{A}}M\rightarrow
B\mathrel{\mathop{\wedge}\limits_{A}}N$ is an equivalence of spectra. Then the
natural map $\lambda^{\wedge}_{f}:M^{\wedge}_{f}\rightarrow N^{\wedge}_{f}$ is
an equivalence of spectra. In particular, if
$B\mathrel{\mathop{\wedge}\limits_{A}}M\simeq*$, then $M^{\wedge}_{f}\simeq*$.
Remark: By the derived smash product, we mean the smash product construction
applied to cofibrant replacements for $M$ and $N$. If they are already
cofibrant, then one can apply the smash product directly.
###### Proposition 2.21
Suppose that $A$ is a (-1)-connected commutative ring spectrum, and that $M$
is a connective $A$-module. Suppose further that we have a homomorphism of
commutative ring spectra $A\rightarrow\mathbb{H}_{l}$, and that for all $n$
and $i$, we have that the groups $Tor^{\pi_{0}A}_{i}(\pi_{n}M,\mathbb{F}_{l})$
vanish. Then we have
$M\mathrel{\mathop{\wedge}\limits_{A}}\mathbb{H}_{l}\simeq*$, and in
particular $M^{\wedge}_{\mathbb{H}_{l}}\simeq*$.
Proof: We have the Künneth spectral sequence converging to
$\pi_{*}(M\mathrel{\mathop{\wedge}\limits_{A}}\mathbb{H}_{l})$, whose
$E_{2}$-term is $Tor_{\pi_{*}A}(\pi_{*}M,\mathbb{F}_{l})$. We also have the
inclusion $\pi_{0}A\hookrightarrow\pi_{*}A$, as the inclusion of the degree
zero component in the graded ring $\pi_{*}A$. We let $B_{*}$ denote the
commutative graded ring
$\mathbb{F}_{l}\mathrel{\mathop{\otimes}\limits_{\pi_{0}A}}\pi_{*}A$. We can
now describe the functor
$\mathbb{F}_{l}\mathrel{\mathop{\otimes}\limits_{\pi_{*}A}}-$ from graded
$\pi_{*}A$-modules to graded $\mathbb{F}_{l}$-vector spaces as the composite
of the functor $\mathbb{F}_{l}\mathrel{\mathop{\otimes}\limits_{\pi_{0}A}}-$
from graded $\pi_{*}A$-modules to graded $B_{*}$-modules with the functor
$\mathbb{F}_{l}\mathrel{\mathop{\otimes}\limits_{B_{*}}}-$ from the category
of graded $B_{*}$-modules to graded $\mathbb{F}_{l}$-vector spaces. The
Grothendieck spectral sequence now shows that the vanishing of the derived
functors of the first functor forces the vanishing of the derived functors of
the composite. $\square$
###### Proposition 2.22
Let $f:A\rightarrow B$ be a homomorphism of $(-1)$-connected commutative ring
spectra, with $\pi_{0}(f)$ surjective, and let $M$ be an $A$-module. Let $I$
denote the kernel of $\pi_{0}(f)$. Suppose that the $\pi_{0}A$-modules
$\pi_{j}M$ are such that for each $j$, there is an integer $e_{j}$ so that
$I^{e_{j}}\cdot\pi_{j}M=\\{0\\}$. Then the natural map $\eta:M\rightarrow
M^{\wedge}_{f}$ is an equivalence.
Proof: Theorem 7.1 of [7] reduces us to the case where $A$ is the ring
$\pi_{0}A$, $B$ is the ring $\pi_{0}B$, and $M$ is a module for which
$I^{e}M=\\{0\\}$ for some $e$. We then obtain the finite filtration
$\\{0\\}=I^{e}M\subseteq I^{e-1}M\subseteq\cdots\subseteq I^{2}M\subseteq
IM\subseteq M$
each of whose quotients is annihilated by $I$. It follows from (6) of
Proposition 3.2 of [7] that
$\eta:I^{s}M/I^{s+1}M\rightarrow(I^{s}M/I^{s+1}M)^{\wedge}_{f}$ is an
equivalence. A straightforward application of Proposition 2.19 now gives the
result. $\square$
The following technical criterion will be useful in later sections.
###### Proposition 2.23
Let
$M_{0}\stackrel{{\scriptstyle
f_{0}}}{{\rightarrow}}M_{1}\stackrel{{\scriptstyle
f_{1}}}{{\rightarrow}}M_{2}\stackrel{{\scriptstyle
f_{2}}}{{\rightarrow}}\cdots$
be a directed system of module spectra over the commutative ring spectrum $A$.
Suppose further that we are given a homomorphism $f:A\rightarrow B$ of
commutative ring spectra. Finally, suppose that the $Ind$-groups
$\pi_{t}(B\mathrel{\mathop{\wedge}\limits_{A}}M_{0})\rightarrow\pi_{t}(B\mathrel{\mathop{\wedge}\limits_{A}}M_{1})\rightarrow\pi_{t}(B\mathrel{\mathop{\wedge}\limits_{A}}M_{2})\rightarrow\cdots$
are $Ind$-trivial. Then the $A$-module $M_{\infty}=\mbox{
}\mathrel{\mathop{hocolim}\limits_{r}}M_{r}$ satisfies
$B\mathrel{\mathop{\wedge}\limits_{A}}M_{\infty}\simeq*$, and therefore
$(M_{\infty})^{\wedge}_{f}\simeq*$.
Proof: Because smash products are colimit constructions, it is possible to
commute the constructions $B\mathrel{\mathop{\wedge}\limits_{A}}-$ and colim,
which shows that $B\mathrel{\mathop{\wedge}\limits_{A}}M_{\infty}\simeq*$.
$\square$
One situation in which one can prove that the hypotheses of Proposition 2.23
are satisfied is treated in the following result.
###### Proposition 2.24
Suppose $f:A\rightarrow B$ is a homomorphism of $(-1)$-connected commutative
ring spectra, which is surjective on $\pi_{0}$, and suppose $L$ is a
connective commutative $A$-algebra. , equipped with a homomorphism of
commutative $A$-algebras $\epsilon_{L}:L\rightarrow B$ so that the diagram
$\begin{diagram}$
commutes. Let $I_{A}\subseteq\pi_{0}A$ denote the ideal $Ker(f)$. Consider now
an Ind-$A$-module of the form
$L\stackrel{{\scriptstyle\times\alpha_{0}}}{{\longrightarrow}}L\stackrel{{\scriptstyle\times\alpha_{1}}}{{\longrightarrow}}L\stackrel{{\scriptstyle\times\alpha_{2}}}{{\longrightarrow}}\cdots$
(2–4)
where $\alpha_{i}$ is an element of the ideal $I_{A}\cdot\pi_{0}L$. Then Ind-
groups
$\pi_{t}(B\mathrel{\mathop{\wedge}\limits_{A}}L)\stackrel{{\scriptstyle\pi_{t}(\times\alpha_{0})}}{{\longrightarrow}}\pi_{t}(B\mathrel{\mathop{\wedge}\limits_{A}}L)\stackrel{{\scriptstyle\pi_{t}(\times\alpha_{1})}}{{\longrightarrow}}\pi_{t}(B\mathrel{\mathop{\wedge}\limits_{A}}L)\stackrel{{\scriptstyle\pi_{t}(\times\alpha_{2})}}{{\longrightarrow}}\cdots$
(2–5)
are Ind-trivial, and therefore by Proposition 2.23, the homotopy colimit
$L_{\infty}=\mbox{}\mathrel{\mathop{hocolim}\limits_{i}}L$ satisfies
$B\mathrel{\mathop{\wedge}\limits_{A}}L_{\infty}\simeq*$.
Proof: From the theory of derived tensor products for modules over ring
spectra, as developed in [33], there is a spectral sequence with $E_{2}$-term
$Tor^{\pi_{*}A}_{*}(\pi_{*}B,\pi_{*}L)$ converging to
$\pi_{*}(B\mathrel{\mathop{\wedge}\limits_{A}}L)$. We now have an inductive
system of spectral sequences corresponding to diagram 2–5 above, and claim
that maps in the system are identically zero. To prove this, we must show that
any $\lambda\in I_{A}\cdot\pi_{0}L$ induces the zero map on all $Tor$-groups.
To see this, we write such a $\lambda$ as
$\lambda=\sum_{s}i_{s}\lambda_{s}$
where $i_{s}\in I_{A}$, and $\lambda_{s}\in\pi_{0}L$. If we can prove that
each $i_{s}$ induces the zero map on $Tor$-groups, we will have verified the
claim. But this is clear, since multiplication by elements of $I_{A}$ induce
the zero map on $\pi_{*}B$. Further, the colimit of the spectral sequences
converges to the colimit of the homotopy groups due to the strong convergence
of the spectral sequences which is guaranteed by the connectivity hypothes on
$A$ and $B$, which gives the result. $\square$
###### Corollary 2.6
Suppose we are given $A,B,L,\mbox{ and }I_{A},$ are as in Proposition 2.24,
and that $\epsilon_{L}:L\rightarrow B$ is a homomorphism of commutative ring
spectra so that the diagram
$\begin{diagram}$
commutes. Let $I_{L}$ be the kernel of $\pi_{0}(\epsilon_{L})$. Suppose
further that there is an integer $N$ so that $I_{L}^{N}\subseteq
I_{A}\cdot\pi_{0}(L)$. We suppose we are given an Ind-$A$-module as in (2–4)
above, but with the assumption that $\alpha_{i}\in I_{L}$ rather than that
$\alpha_{i}\in I_{A}$. Then the Ind-groups in (2–5) are trivial, and we again
have that $B\mathrel{\mathop{\wedge}\limits_{A}}L_{\infty}\simeq*$
Proof: We note that any $N$-fold composite in the system (2–5) has its image
in $I_{A}\cdot\pi_{0}(L)$, and a standard cofinality argument gives the
result. $\square$
## 3 Totally torsion free groups
We recall that a topological group $G$ is said to be profinite if it is
Hausdorff, compact, and totally disconnected. These properties imply that as a
topological group, it can be expressed as an inverse limit of finite groups.
Let $l$ denote a prime. A profinite $l$-group is a profinite group whose only
finite quotients are finite $l$-groups. For any profinite group $G$, we define
the commutator subgroup of $G$, $[G,G]$ to be the closure of the subgroup
generated by elements of the form $ghg^{-1}h^{-1}$ in $G$. The subgroup
$[G,G]$ is normal, and the quotient $G/[G,G]$ will be denoted by $G^{ab}$. It
is the universal profinite abelian quotient of $G$. One also has the
following.
###### Proposition 3.1
Any closed subgroup of a profinite group is itself a profinite group. The
quotient of a profinite group by a closed normal subgroup is also a profinite
group. In particular, any subgroup containing a closed subgroup of finite
index is itself a closed subgroup.
This is in $\S$2.1 of [32].
We say a profinite $l$-group $G$ is totally torsion free if for every closed
subgroup $K\subseteq G$ the quotient group $K^{ab}$ is torsion free. Free
abelian pro-$l$ groups and free pro-$l$ groups are totally torsion free, and
free products (in the category of profinite pro-$l$ groups) of totally torsion
free pro-$l$ groups are totally torsion free. However, the condition is quite
restrictive, as the following example shows.
Example: The $2$-adic Heisenberg group, i.e. the group of $3\times 3$ upper
triangular matrices with entries in $\mathbb{Z}_{2}$ and ones on the diagonal
is not TTF, although it is easily seen to be torsion free. To see this, we
note that the subgroup of matrices of the form
$\left[\begin{array}[]{ccc}1&2k&l\\\ 0&1&m\\\ 0&0&1\end{array}\right]$
has abelianization isomorphic to
$\mathbb{Z}_{2}\oplus\mathbb{Z}_{2}\oplus\mathbb{Z}/2\mathbb{Z}$
The matrix
$\left[\begin{array}[]{ccc}1&0&1\\\ 0&1&0\\\ 0&0&1\end{array}\right]$
is not a commutator, but its square is.
###### Proposition 3.2
Let $k$ be a field containing all roots of unity, with absolute Galois group
$G_{k}$. Then $G_{k}$ is totally torsion free.
Proof: It clearly suffices to prove that $G_{k}^{ab}$ is torsion free for all
fields $k$ containing all the roots of unity. Kummer theory (see e.g. [26],
§VI.8) tells us that $G_{k}^{ab}$ is the Galois group of the infinite
extension $k^{ab}=\bigcup_{m}k((k^{*})^{\frac{1}{m}})$, and therefore that any
non-identity element $g\in G_{k}^{ab}$ acts non-trivially on an element of the
form $\sqrt[m]{\kappa}$ for some $m$ and $\kappa\in k$. Suppose now that
$g^{n}=e$ in $G_{k}^{ab}$. We have
$g\cdot\sqrt[m]{\kappa}=\zeta\sqrt[m]{\kappa}$ for some $\zeta\in\mu_{m}(k)$.
Now choose any $n$-th root $\xi$ of $\sqrt[m]{\kappa}$. We see that
$g\cdot\xi=\eta\xi$ for some $mn$-th root of unity $\eta$ satisfying
$\eta^{n}=\zeta$. Now, $g^{n}\cdot\xi=\eta^{n}\xi=\zeta\xi\neq\xi$, which
contradicts the assumption that $g$ is $n$-torsion. $\square$
Remark: This structural property of Galois groups will allow us to identify
completion constructions based on representation theory with $K$-spectra of
certain schemes which appear to be the algebraic geometric version of the
classifying space construction in topology. It would be interesting to
consider how this property relates with the Bloch-Kato conjecture.
## 4 Construction of some continuous affine $G$-schemes
Let $G$ be a profinite $l$-group. We will need to construct continuous affine
$G$-actions so as to be able to assemble the scheme ${\cal E}G$ we require. In
this section, we will construct such actions attached to affine
representations over $\mathbb{Z}_{l}$ of $G$. By taking (infinite) products of
schemes constructed in this way, we will obtain the scheme ${\cal E}G$
discussed in the introduction when $G$ is totally torsion free.
###### Definition 4.1
An affine $\mathbb{Z}_{l}$-module is a pair $(F,X)$, where $F$ is a free
finitely generated $\mathbb{Z}_{l}$-module, and $X$ is a continuous simply
transitive $F$-space. Here $F$ is regarded as a topological abelian group, and
the simple transitivity means that for every $x\in X$, the natural map
$F\rightarrow X$ given by $f\rightarrow f\cdot x$ is a homeomorphism. A based
affine $\mathbb{Z}_{l}$-module is a triple $(F,X,B)$ so that the pair $(F,X)$
is an affine $\mathbb{Z}_{l}$-module and $B$ is a basis for the
$\mathbb{Z}_{l}$-module $F$. An isomorphism of affine $\mathbb{Z}_{l}$-modules
from $(F,X)$ to $(F^{\prime},X^{\prime})$ is a pair $(\varphi,\chi)$, where
$\varphi:F\rightarrow F^{\prime}$ is a continuous isomorphism of
$\mathbb{Z}_{l}$-modules and $\chi:X\rightarrow X^{\prime}$ is a
homeomorphism, and where the relationship
$\chi(f\cdot x)=\varphi(f)\cdot\chi(x)$
holds for all $f$ and $x$. An isomorphism $\varphi$ of based affine
$\mathbb{Z}_{l}$-modules is an isomorphism of the underlying affine
$\mathbb{Z}_{l}$-modules which preserves $\mathbb{Z}_{l}$-multiples of basis
elements, i.e. for all $z\in\mathbb{Z}_{l}$ and $b\in B$ there exist
$z^{\prime}\in\mathbb{Z_{l}}$ and $b^{\prime}\in B$ so that
$\varphi({zb})=z^{\prime}b^{\prime}$. If for all $z$ and $b$, there is an
element $z^{\prime}\in\mathbb{Z}_{l}$ so that $\varphi(zb)=z^{\prime}b$, then
the isomorphism is said to be a translation.
###### Proposition 4.1
For $F$ a free, finitely generated $\mathbb{Z}_{l}$-module and a simply
transitive $G$-space $X$, we have an isomorphism
$\mbox{Aut}(F,X)\cong GL(F)\ltimes F$
For a based affine $\mathbb{Z}_{l}$-module $(F,X,B)$ we have
$\mbox{Aut}(F,X,B)\cong\Sigma_{B}\ltimes F$
where $\Sigma_{B}$ denotes the symmetric group of permutations of the set $B$.
Proof: Clear from the definitions of the morphisms. $\square$
Remark: The group $\mbox{Aut}(F,X,B)$ is a profinite group. For every $n>0$,
we have the quotient of $X$ by the equivalence relation $\simeq_{n}$ given by
$x\simeq_{n}y$ if and only if there is an element $f\in l^{n}F$ so that
$f\cdot x=y$. These equivalence relations are respected by the action of
$\mbox{Aut}(F,X,B)$, and the quotients are finite sets.
Let $(F,X,B)$ be a based affine $\mathbb{Z}_{l}$-module, and let
$\mathfrak{A}=\mbox{Aut}(F,X,B)$. Also, let
$\mathfrak{T}\subseteq\mathfrak{A}$ denote the subgroups of translations. We
now construct a family $\\{\rho_{s}\\}$ of $k$-linear representations of
$\mathfrak{A}$. Let $T^{\infty}$ denote the inverse limit
$\mathrel{\mathop{lim}\limits_{\leftarrow}}\mu_{l^{n}}$, and choose a
continuous identification
$\alpha:\mathbb{Z}_{l}\stackrel{{\scriptstyle\sim}}{{\rightarrow}}T^{\infty}$.
Let $\pi_{s}:T^{\infty}\rightarrow\mu_{l^{s}}$ be the projection. Fixing an
element $b\in B$, we define a subgroup
$\mathfrak{A}^{\prime}\subseteq\mathfrak{A}$ to be $p^{-1}(\Sigma_{B}(b))$,
where $p:\mathfrak{A}\rightarrow\Sigma_{B}$ is the natural projection, and
where $\Sigma_{B}(b)\subseteq\Sigma_{B}$ denotes the stabilizer of the element
$b$. We define a one dimensional character $\sigma_{s}$ of
$\mathfrak{A}^{\prime}$ by setting $\sigma_{s}(b)=\pi_{s}\alpha(1)$,
$\sigma_{s}(b^{\prime})=1$ for $b^{\prime}\neq b$, and
$\sigma_{s}(\Sigma_{B}(b))=\\{1\\}$. We define $\rho_{s}$ to be the induced
representation
$\mbox{Ind}_{\mathfrak{A}^{\prime}}^{\mathfrak{A}}(\sigma_{s})$
Of course, each $\rho_{s}$ now gives an action of $\mathfrak{A}$ on the
$n$-dimensional affine space over $k$ by algebraic automorphisms. We denote
this variety with $\mathfrak{A}$-action by $\mathbb{A}^{n}(k,\rho_{s})$. Note
that the affine space also is equipped with a fixed choice of coordinates,
since it is an induced representation. The coordinates correspond to the
cosets $\Sigma_{B}/\Sigma_{B}(b)$, which in turn correspond to the basis $B$
of $F$. We define a self map
$\theta:\mathbb{A}^{n}(k)\rightarrow\mathbb{A}^{n}(k)$ by requiring that
$\theta(\sum x_{b}b)=\sum x_{b}^{l}b$
###### Proposition 4.2
The map $\theta$ is $\mathfrak{A}$-equivariant as a map from
$\mathbb{A}^{n}(k,\rho_{s+1})$ to $\mathbb{A}^{n}(k,\rho_{s})$
Proof: Immediate. $\square$
We now consider the pro-scheme
$\cdots\stackrel{{\scriptstyle\theta}}{{\rightarrow}}\mathbb{A}^{n}(k,\rho_{s+1})\stackrel{{\scriptstyle\theta}}{{\rightarrow}}\mathbb{A}^{n}(k,\rho_{s})\stackrel{{\scriptstyle\theta}}{{\rightarrow}}\mathbb{A}^{n}(k,\rho_{s-1})\rightarrow\cdots\stackrel{{\scriptstyle\theta}}{{\rightarrow}}\mathbb{A}^{n}(k,\rho_{1})\stackrel{{\scriptstyle\theta}}{{\rightarrow}}\mathbb{A}^{n}(k,\rho_{0})$
It is equipped with the action by $\mathfrak{A}$, and we will denote its
inverse limit by ${\cal A}_{(F,X,B)}$. This inverse limit exists because the
system consists of affine varieties, and is itself the affine scheme
associated to the ring $\bigcup_{n}k[t^{\frac{1}{l^{n}}}]$. The
$\mathfrak{A}$-action is a continuous affine $\mathfrak{A}$-action. We also
construct the related continuous affine $\mathfrak{A}$-action ${\cal
T}_{(F,X,B)}$, which is obtained by deleting the coordinate hyperplanes from
each $\mathbb{A}^{n}(k,\rho_{s})$ for each $s$. We also define ${\cal
A}^{s}_{(F,X,B)}$ and ${\cal T}^{s}_{(F,X,B)}$ to be the representation
$\mathbb{A}^{n}(k,\rho_{s})$ and the result of removing the coordinate
hyperplanes from $\mathbb{A}^{n}(k,\rho_{s})$, respectively.
Remark: The construction of ${\cal A}_{(F,X,B)}$ and ${\cal T}_{(F,X,B)}$
depends on the choice of $b\in B$. It is easy to verify that varying this
choice does not affect ${\cal A}_{(F,X,B)}$ or ${\cal T}_{(F,X,B)}$ up to
isomorphism.
###### Proposition 4.3
The restriction of the $\mathfrak{A}$ action on ${\cal T}_{(F,X,B)}$ to
$\mathfrak{T}$ is an etale $\mathfrak{T}$-action.
Proof: Consider the direct system of rings
$k[t_{1}^{\pm 1},\ldots,t_{n}^{\pm 1}]\subseteq
k[t_{1}^{\pm\frac{1}{l}},\ldots,t_{n}^{\pm\frac{1}{l}}]\subseteq
k[t_{1}^{\pm\frac{1}{l^{2}}},\ldots,t_{n}^{\pm\frac{1}{l^{2}}}]\subseteq\cdots$
where $n$ is the rank of $F$, and let ${\cal L}_{n}$ denote its colimit. This
system is identified with the system of rings one obtains by applying the
affine coordinate ring functor to the system defining ${\cal T}_{(F,X,B)}$.
Further, ${\cal L}_{n}$ is an infinite Galois extension of $k[t_{1}^{\pm
1},\ldots,t_{n}^{\pm 1}]$, with Galois group $\Gamma=\mathbb{Z}_{l}^{n}$. The
restriction of the $\mathfrak{A}$-action on ${\cal T}_{(F,X,B)}$ constructed
above using the linear representations $\rho_{s}$, when restricted to
$\mathfrak{T}$, clearly provide an identification of $\mathfrak{T}\cong F$
with $\Gamma$, and the result now follows. $\square$
###### Definition 4.2
For any pro-$l$ group $G$, an affine $l$-adic representation (or $l$-BAR) is a
continuous homomorphism from $G$ to the automorphism group of a based affine
$\mathbb{Z}_{l}$-module. For any affine $l$-adic representation $\eta$ of $G$
on an affine $\mathbb{Z}_{l}$ module $(F,X,B)$, we obtain actions of $G$ on
the schemes ${\cal A}_{(F,X,B)}$ and ${\cal T}_{(F,X,B)}$ by pulling back the
action of the automorphism group of $(F,X,B)$ on these schemes. We will denote
these $G$-schemes by ${\cal A}_{\eta}$ and ${\cal T}_{\eta}$. We also write
${\cal A}^{s}_{\eta}$ and ${\cal T}_{\eta}^{s}$ for the pullback along $\eta$
of the actions ${\cal A}_{(F,X,B)}^{s}$ and ${\cal T}_{(F,X,B)}^{s}$.
There is a useful description of ${\cal T}_{\eta}^{s}$ which we will need.
###### Proposition 4.4
The scheme ${\cal T}_{\eta}^{s}$ can be identified with the quotient ${\cal
T}_{\eta}/l^{s}\cdot\mathfrak{T}$, where $\mathfrak{T}$ denotes the full
translation subgroup of $\mathfrak{A}$. $G$ acts on this quotient, since the
image of $G$ under $\eta$ normalizes all the subgroups
$l^{s}\cdot\mathfrak{T}$.
Proof: This is just the ring theoretic result that
$\bigcup_{n}k[t_{1}^{\pm\frac{1}{l^{n}}},\ldots,t_{k}^{\pm\frac{1}{l^{n}}}]^{l^{s}\cdot\mathfrak{T}}=k[t_{1}^{\pm\frac{1}{l^{s}}},\ldots,t_{k}^{\pm\frac{1}{l^{s}}}]$
$\square$
###### Proposition 4.5
Let $\eta$ be any $l$-BAR. Then the continuous affine $G$-actions ${\cal
A}_{\eta}\times Spec(L)$ and ${\cal T}_{\eta}\times Spec(L)$ are both
coherent, where $L$ is any field equipped with a continuous action by $G$. As
always, it is assumed that $L$ is equipped with the discrete topology.
Proof: It suffices to prove that the rings
$\bigcup_{n}L[x_{1}^{\frac{1}{l^{n}}},\ldots,x_{s}^{\frac{1}{l^{n}}}]$
and
$\bigcup_{n}L[x_{1}^{\pm\frac{1}{l^{n}}},\ldots,x_{s}^{\pm\frac{1}{l^{n}}}]$
are coherent rings for any field $L$. This follows immediately from
Proposition 2.3 since the inclusions
$L[x_{1}^{\frac{1}{l^{n}}},\ldots,x_{s}^{\frac{1}{l^{n}}}]\hookrightarrow
L[x_{1}^{\frac{1}{l^{n+1}}},\ldots,x_{s}^{\frac{1}{l^{n+1}}}]$
and
$L[x_{1}^{\pm\frac{1}{l^{n}}},\ldots,x_{s}^{\pm\frac{1}{l^{n}}}]\hookrightarrow
L[x_{1}^{\pm\frac{1}{l^{n+1}}},\ldots,x_{s}^{\pm\frac{1}{l^{n+1}}}]$
are all flat. $\square$
Given a possibly infinite set of $l$-BAR’s $\\{\eta_{i}\\}$, the schemes
$\prod{\cal A}_{\eta_{i}}$ and $\prod{\cal T}_{\eta_{i}}$ both are also
continuous affine $G$-schemes.
###### Corollary 4.1
Both $\prod{\cal A}_{\eta_{i}}$ and $\prod{\cal T}_{\eta_{i}}$ are coherent.
Proof: It will suffice to prove that the rings
$\bigcup_{n}L[x_{1}^{\frac{1}{l^{n}}},\ldots,x_{s}^{\frac{1}{l^{n}}},\ldots]$
and
$\bigcup_{n}L[x_{1}^{\pm\frac{1}{l^{n}}},\ldots,x_{s}^{\pm\frac{1}{l^{n}}},\ldots]$
are coherent. This follows easily from Proposition 2.3 since the inclusions
$\bigcup_{n}L[x_{1}^{\frac{1}{l^{n}}},\ldots,x_{s}^{\frac{1}{l^{n}}}]\hookrightarrow\bigcup_{n}L[x_{1}^{\frac{1}{l^{n}}},\ldots,x_{s+1}^{\frac{1}{l^{n}}}]$
and
$\bigcup_{n}L[x_{1}^{\pm\frac{1}{l^{n}}},\ldots,x_{s}^{\pm\frac{1}{l^{n}}}]\hookrightarrow\bigcup_{n}L[x_{1}^{\pm\frac{1}{l^{n}}},\ldots,x_{s+1}^{\pm\frac{1}{l^{n}}}]$
are all flat. $\square$
We now have the following definition.
###### Definition 4.3
Let $G$ be a profinite group, and suppose $K$ is any closed subgroup of finite
index. We say an affine $l$-adic representation $\eta$ of $G$ is $K$-principal
if the stabilizer in $G$ of any $x\in X_{\eta}$ is equal to $K$. Note that it
follows that $K$ is normal. Clearly, the condition is equivalent to the
condition that $X$ is a free $G/K$-set.
The following lemma is an $l$-adic analogue of the results of Auslander and
Kuranishi concerning the structure of fundamental groups of flat manifolds
[3].
###### Lemma 4.1
Let $K\subseteq G$ be a closed normal subgroup of a profinite group $G$. Then
a based affine $l$-adic representation $\eta$ of a $G$ group is $K$-principal
if and only if the image of $\eta$ in the automorphism group of $X_{\eta}$ is
torsion free.
Proof: We let $F$ be a free finitely generated $\mathbb{Z}_{l}$-module with
basis $B$, so that
$(F_{\eta},X_{\eta},B_{\eta})\cong(F,F,B)$
where $F$ acts on itself by translation. The automorphism group of $(F,F,B)$
is a semidirect product $\Sigma_{B}\ltimes F$. The group $\Sigma_{B}$ acts on
$F$ by permuting the basis $B$, and we describe the action of an element
$(\sigma,f)$, where $\sigma\in\Sigma_{B}$ and $f\in F$ by $(\sigma,f)\cdot
f^{\prime}=\sigma\cdot f^{\prime}+f$. Given an element $(\sigma,f)$ of
$\Sigma_{B}\ltimes F$, the action of $(\sigma,f)$ on $F$ has a fixed point if
and only if $f=v-\sigma v$ for some $v\in F$. This condition is in turn
equivalent to the requirement that any $\mathbb{Z}_{l}$-linear function which
satisfies $f\cdot\sigma=f$ vanishes on $f$. Functions satisfying this
condition form a free $\mathbb{Z}_{l}$-module of rank equal to the number of
cycles in the cycle decomposition of the permutation $\sigma$, with the
characteristic functions on orbits forming a basis. On the other hand,
$(\sigma,f)$ has finite order $n$ if and only if $\sigma^{n}=e$ and
$f+\sigma\cdot f+\sigma^{2}\cdot f+\cdots+\sigma^{n-2}\cdot
f+\sigma^{n-1}\cdot f=0$
Let $\\{\mathfrak{o}_{i}\\}_{i}$ denote the collections of orbits under action
of $\sigma$, and let $\pi_{i}$ denote the projection on the set of basis
elements in $\mathfrak{o}_{i}$. The summands corresponding to the orbits are
invariant under $\sigma$, and we obtain a $\sigma$-invariant direct sum
decomposition
$F\cong\bigoplus_{i}F_{i}$
where $F_{i}$ is the span of the basis elements belonging to
$\mathfrak{o}_{i}$. The above equation now demonstrates that
$\sum_{s=0}^{n-1}\sigma^{s}(\pi_{i}(f))=0$
It follows that all the $\sigma$-invariant $\mathbb{Z}_{l}$-linear functions
on $F$ vanish on $f$. So we have proved that for any element $(\sigma,f)\in
Aut(F,F,B)$, the condition that $(\sigma,f)$ as a fixed point under its action
on $F$ is equivalent to its having finite order. The result now follows.
$\square$
We will write ${\cal T}_{\eta}/G$ for the orbit scheme of ${\cal T}_{\eta}$.
We now have the following result.
###### Proposition 4.6
Let $G$ be a profinite group, with $N\subseteq G$ a closed normal subgroup.
Let $\eta$ denote a based affine $l$-adic representation of $G$. If $\eta$ is
$N$-principal, then ${\cal T}_{\eta}$ is an etale $G/N$-scheme.
Proof: Let $(F_{\eta},X_{\eta},B_{\eta})$ be the based $l$-adic representation
space, so $\eta$ is a continuous homomorphism $G\rightarrow\mathfrak{A}$,
where $\mathfrak{A}$ denotes the automorphism group of
$(F_{\eta},X_{\eta},B_{\eta})$. As before, let $\mathfrak{T}$ denote the
translation subgroup of $\mathfrak{A}$. The subgroup $\eta^{-1}(\mathfrak{T})$
is now a closed subgroup of finite index in $G$, containing $N$, and therefore
projects to a closed subgroup $T$ of finite index in $G/K$. Let $\mathfrak{X}$
denote the set of all morphisms from schemes of the form $Spec(F)$, where $F$
is an algebraically closed field, to the inverse limit ${\cal T}_{\eta}$. In
order to prove the freeness of the $G/N$-scheme, we must prove the freeness of
he $G/N$ action on the set $\mathfrak{X}$. We know from Proposition 4.3 that
the restriction of the action on $\mathfrak{X}$ to the subgroup of finite
index $T$ is free, and that the group $G/N$ is torsion free. In order to prove
the freeness of the full $G/N$-action, we suppose that there is are elements
$\gamma\in G/N$, with $\gamma\neq e$, and $x\in\mathfrak{X}$, so that
$\gamma\cdot x=x$. It would follow that for any $n$, the element $\gamma^{n}$
would also fix $x$. But for sufficiently large $n$, $\gamma^{n}\in T$, and
further that $\gamma^{n}\neq e$, because of the torsion freeness proved in
Proposition 4.1. The freeness of the $N$ -action now contradicts the existence
of $\gamma$. $\square$
We actually need a slightly stronger result, namely that for a fixed
$N$-principal $l$-BAR, there is an $s\geq 0$ so that the scheme ${\cal
T}^{s}_{\eta}$ is acted on freely by the quotient group
$G/N\cdot\eta^{-1}(l^{s^{\prime}}\mathfrak{T})$ for all $s^{\prime}\geq s$.
This will require a few algebraic facts.
Consider any extension of groups
$0\rightarrow A\rightarrow H\rightarrow Q\rightarrow 0$
where $Q$ is a finite cyclic group of order $l$, and $A$ is abelian. The
extension gives an action of $Q$ on $A$ by automorphisms. Let $N(A)\subseteq
A$ denote the subgroup of elements of the form $a+\sigma\cdot
a+\sigma^{2}\cdot a+\cdots\sigma^{p-1}a$, where $\sigma\in Q$ is a generator.
Of course the group is independent of the choice of generator. We define a
function $c:Q\rightarrow A/N(A)$ as follows. For any element $\sigma\in Q$,
select a lift $\overline{\sigma}\in H$, i.e. so that
$\pi\overline{\sigma}=\sigma$, and define $c(\sigma)=\overline{\sigma}^{l}$.
It is easy to check that $c(\sigma)$ is independent of the choice of
$\overline{\sigma}$.
###### Lemma 4.2
An element $\sigma\in Q$ lifts to an element of order $l$ if and only if
$c(\sigma)=0$.
Proof: The “only if” part is clear, since if $\sigma$ lifts to an element of
order $l$, $c(\sigma)$ is by definition equal to zero. Suppose on the other
hand that $c(\sigma)=0$, which means that there is a lift $\overline{\sigma}$
and an element $\alpha\in A$ so that
$\overline{\sigma}^{l}=\alpha+\sigma\alpha+\cdots\sigma^{l-1}\alpha$
If we now set $\overline{\sigma}^{\prime}=\overline{\sigma}\cdot(-\alpha)$,
then $(\overline{\sigma}^{\prime})^{l}=e$. $\square$
###### Lemma 4.3
Let $\Gamma$ denote a pro-$l$ group which fits in an exact sequence
$\\{e\\}\rightarrow
L\rightarrow\Gamma\stackrel{{\scriptstyle\pi}}{{\rightarrow}}Q\rightarrow\\{e\\}$
where $Q$ is a finite $l$-group, and $L\cong\mathbb{Z}_{l}^{n}$. Let
$\Gamma^{\prime}\subseteq\Gamma$ be a subgroup so that the composite
$\Gamma^{\prime}\rightarrow\Gamma\rightarrow Q$
is surjective. Suppose further that $\Gamma^{\prime}$ is torsion free, and
that we are given a continuous $\Gamma$-set $X$ on which $\Gamma^{\prime}$ and
$L$ both act freely. Then there is an $s\geq 0$ so that
$\Gamma^{\prime}/\Gamma^{\prime}\cap l^{s^{\prime}}\cdot L$ acts freely on the
set $X/l^{s^{\prime}}\cdot L$ for all $s^{\prime}\geq s$. The number $s$
depends only on the groups involved, not the $\Gamma$-set $X$.
Proof: We first claim that it suffices to prove the result for $Q$ a cyclic
finite $l$-group. To see this, we suppose we are able to prove the result for
each subgroup
$\Gamma^{\prime}_{C}=(\pi|\Gamma^{\prime})^{-1}(C)\subseteq\Gamma^{\prime}$.
This means that for every cyclic subgroup $C\subseteq Q$, there is an
$s(C)\geq 0$ so that $\Gamma_{C}^{\prime}/\Gamma_{C}^{\prime}\cap
l^{s^{\prime}}L$ acts freely on $X/l^{s^{\prime}}L$ for all $s^{\prime}\geq
s(C)$ . Since there are only finitely many cyclic subgroups of $Q$, we may set
$S=\mathrel{\mathop{max}\limits_{C}}s(C)$. It now follows that
$\Gamma^{\prime}/\Gamma^{\prime}\cap l^{s^{\prime}}L$ acts freely on
$X/l^{s^{\prime}}L$ whenever $s^{\prime}\geq S$, since all the groups
$\Gamma_{C}^{\prime}/\Gamma_{C}^{\prime}\cap l^{s^{\prime}}L$ do, and since
these the union of these subgroups is all of
$\Gamma^{\prime}/\Gamma^{\prime}\cap l^{s^{\prime}}L$. From this point on, we
assume that $Q$ is cyclic of prime power order. The next reduction is that it
will suffice to prove the result for $Q$ a cyclic group of order $l$. To see
this, we first assume that the result holds for all groups $\Gamma^{\prime}$
with $Q$ of order $l$, and let $\Gamma^{\prime\prime}\subseteq\Gamma^{\prime}$
denote the subgroup $\pi^{-1}(Q_{l})$, where $Q_{l}$ is the unique subgroup of
order $l$. Suppose $\gamma\in\Gamma^{\prime}-\Gamma^{\prime\prime}$, and let
$l^{k}$ denote the order $\pi(\gamma)\in Q$. We then have that
$\gamma^{l^{k-1}}\in Q_{l}$, and $\gamma^{l^{k-1}}\neq e$. By our assumption,
there is an $s\geq 0$ so that the group
$\Gamma^{\prime\prime}/\Gamma^{\prime\prime}\cap l^{s}L$ acts freely on
$X/l^{s^{\prime}}\cdot LX$ for all $s^{\prime}\geq s$. This means that
$\gamma$ also acts freely, since a power of it acts freely. This give the
reduction to the case where $Q$ has order $l$. To handle this case, we note
that we have the function $c:Q\rightarrow L\cap\Gamma^{\prime}$ defined above.
By Lemma 4.2 above, the torsion freeness of $\Gamma^{\prime}$ means that
$c(\sigma)\neq 0$ in the group $L\cap\Gamma^{\prime}/N(L)\cap\Gamma^{\prime}$.
But it is clear that $L\cap\Gamma^{\prime}/N(L)\cap\Gamma^{\prime}$ is
identified with the inverse limit
$\mathrel{\mathop{lim}\limits_{\mathrel{\mathop{\longleftarrow}\limits_{s}}}}L\cap\Gamma^{\prime}/(N(L)\cap\Gamma^{\prime}+l^{s}L\cap\Gamma^{\prime})$
so for sufficiently large $s$, we have $\overline{\sigma}^{l}\neq 0$ in
$L\cap\Gamma^{\prime}/l^{s}L\cap\Gamma^{\prime}$, for every possible lift
$\overline{\sigma}$ of $\sigma$. Since the $L$ action on $X$ is assumed to be
free, it follows that $L\cap\Gamma^{\prime}/l^{s}L\cap\Gamma^{\prime}$ acts
freely on $X/l^{s}\cdot L$ for sufficiently large $s$. We now have that every
lift of $\sigma$ acts freely for sufficiently large $s$. It is clear also that
the same calculation applies to all non-identity powers of $\sigma$, which
together with the fact that $L/l^{s}\cdot L$ acts freely on $X/l^{s}\cdot L$,
as always for sufficiently large $s$, gives the result. $\square$
###### Proposition 4.7
Let $\eta$ be a $N$-principal $l$-BAR for some closed normal subgroup of $G$.
Then there exists an $s\geq 0$ so that for every $s^{\prime}\geq s$, the
action of the finite group $G/(N\cdot\eta^{-1}(l^{s^{\prime}}\mathfrak{T}))$
on ${\cal T}^{s^{\prime}}_{\eta}$ is etale.
Proof: The $G/N$-action on ${\cal T}_{\eta}$ is etale. Let $\Omega$ denote an
algebraically closed field and let $i:Spec\Omega\rightarrow{\cal T}_{\eta}/G$
be a morphism. Then we have the set of all maps $\Omega\rightarrow{\cal
T}_{\eta}$ for which the diagram
$\begin{diagram}$
commutes, which we have denoted by $\mathfrak{H}(A[{\cal
T}_{\eta}],\Omega,i)$. The etaleness of the $G/N$-action on ${\cal T}_{\eta}$
means that the $G/N$-action on $\mathfrak{H}(A[{\cal T}_{\eta}],\Omega,i)$ is
free. The group $G/N$ is identified by the representation with a subgroup of
$\Sigma_{n}(l)\ltimes\mathbb{Z}_{l}^{n}$, where $\Sigma_{n}(l)$ denotes an
$l$-Sylow subgroups,and is torsion free, so satisfies the hypotheses of Lemma
4.3 above. Note that the hypothesis that $L$ acts freely in Lemma 4.3 is
satisfied due to Proposition 4.3. The group $\Gamma$ will be
$I\ltimes\mathbb{Z}_{l}^{n}$, where $I$ is the image of the composite
$G/K\rightarrow\Sigma_{n}(l)\times\mathbb{Z}_{l}^{n}\rightarrow\Sigma_{n}(l)$
and the affine representation $\eta$ identifies $G/N$ with a subgroup of
$\Gamma$. It now follows that there is an $s$ so that for all $s^{\prime}\geq
s$, the group $G/(N\cdot\eta^{-1}(l^{s^{\prime}}\mathfrak{T}))$ acts freely on
the orbit set
$\mathfrak{H}(A[{\cal T}_{\eta}],\Omega,i)/(l^{s^{\prime}}\mathfrak{T})$
But, because of the description of ${\cal T}_{\eta}^{s}$ given in Proposition
4.4, this orbit set is just the set $\mathfrak{H}(A[{\cal
T}_{\eta}^{s}],\Omega,i)$, which gives the result. Note that we are using the
self evident fact that if a group $G$ acts freely on a set $X$, and $N$ is a
normal subgroup of $G$, then the naturally defined action of $G/N$ on $X/N$ is
also free. $\square$
We record a useful property of $K$-principal actions.
###### Proposition 4.8
Let $\eta_{1}$ and $\eta_{2}$ be affine $l$-adic representations of a
profinite group $G$. Then we have the evident product affine $l$-adic
representation $\eta_{1}\times\eta_{2}$. Let $\\{e\\}\subseteq
K^{\prime}\subseteq K\subseteq G$ be closed subgroups. Suppose that $\eta_{1}$
is $K$-principal, and that $\eta_{2}|K$ is $K^{\prime}$-principal. Then
$\eta_{1}\times\eta_{2}$ is $K^{\prime}$-principal. In particular, if
$\eta_{i}$ is $K_{i}$-principal, for $K_{i}\subseteq G$, then
$\eta_{1}\times\eta_{2}$ is $K_{1}\cap K_{2}$-principal.
Proof: Clear. $\square$
We also consider the notion of induction of affine based $l$-adic
representations. Let $\eta=(F_{\eta},X_{\eta},B_{\eta})$ be a based affine
$l$-adic representation of $K$, where $K$ is a closed subgroup of finite index
in a profinite group $G$. Then we define an induced based affine $l$-adic
representation
$i_{K}^{G}(\eta)=(F_{i_{K}^{G}(\eta)},X_{i_{K}^{G}(\eta)},B_{i_{K}^{G}(\eta)})$
as follows. The sets $F_{i_{K}^{G}(\eta)}$ and $X_{i_{K}^{G}(\eta)}$ are
defined as the sets of equivariant functions $F^{K}(G,F_{\eta})$ and
$F^{K}(G,X_{\eta})$, respectively, where $K$ acts by left multiplication on
$G$ and via $\eta$ on $F_{\eta}$ and $X_{\eta}$. The group $G$ acts by right
multiplication on $G$, yielding a left action on the two sets. The free action
of $F^{K}(G,F_{\eta})$ on $F^{K}(G,X_{\eta})$ is given by pointwise
multiplication. To define $B_{i_{K}^{G}(\eta)}$, we first define elements
$f_{b}\in F_{i_{K}^{G}(\eta)}$ for each $b\in B_{\eta}$ by setting
$f_{b}(e)=b$. The remainder of $B_{i_{K}^{G}(\eta)}$ is obtained by
translating the elements using the $G$ action. It is clear that this set will
give a basis for $F_{i_{K}^{G}(\eta)}$.
We say an affine based $l$-adic representation $\eta$ of a profinite group $G$
is a translation representation if for every $g\in G$, the automorphism
$\eta(g)$ is a translation.
###### Proposition 4.9
Let $G$ be a profinite group, and let $K\lhd G$ be a normal subgroup of finite
index. Suppose we are given a one-dimensional based affine $l$-adic
representation $\eta$ of $K$. Then we have the following.
1. 1.
The restriction of $i_{K}^{G}(\eta)$ to $K$ is a translation representation.
2. 2.
The restriction of $i_{K}^{G}(\eta)$ to $K$ is $K^{\prime}$-principal, where
$K^{\prime}$ is defined as follows. For each element $\gamma\in G/K$, we
denote by $\eta^{\gamma}$ the result of composing $\eta$ with the conjugation
action of a coset representative of $\gamma$. Then $K^{\prime}$ is given by
$K^{\prime}=\bigcap_{\gamma\in G/K}Ker(\eta^{\gamma})$
## 5 Construction of ${\cal E}G$
In this section, we will use products of the actions constructed in the last
section to build pro-objects in the category of continuous $G$-actions of a
profinite $l$-group which will satisfy the requirements formulated in the
introduction. Throughout this section, we will allow ourselves to refer to a
based affine $l$-adic representation as an $l$-BAR, in the interest of
brevity. We begin with a technical lemma.
###### Lemma 5.1
Let $G$ be a totally torsion free pro-$l$-group. Let $G_{0}\subseteq G$ be any
closed subgroup of finite index. Then there exists a closed normal subgroup
$K\subseteq G_{0}$ and a $K$-principal $l$-BAR $\eta$ of $G$.
Proof: We first observe that it suffices to prove the result for $G_{0}$
normal in $G$, since there is a always a normal subgroup of finite index
contained in $G_{0}$.
To construct an $l$-BAR which is $K$-principal for some $K\subseteq G_{0}$, we
first note that the set of subgroups of $G$ containing $G_{0}$ consists
entirely of closed subgroups, by Proposition 3.1. We now assume by induction
that we are able to construct such $l$-BARs for all closed subgroups of index
less than the index of $G_{0}$ in $G$. We construct such $l$-BAR’s for every
normal subgroup containing $G_{0}$ properly. There are obviously only finitely
many such subgroups. We consider the finite list $\\{N_{i}\\}_{i=1}^{s}$ of
all normal subgroups of $G$ which contain $G_{0}$ properly (they are
necessarily closed), and for each one we construct a $K_{i}$-principal $l$-BAR
$\eta_{i}$ of $G$, where $K_{i}\subseteq N_{i}$. The $l$-BAR
$\eta_{1}\times\eta_{2}\times\cdots\times\eta_{s}$ is now $K_{1}\cap\cdots\cap
K_{s}$-principal by repeated application of Proposition 4.8. We now consider
the subgroup $G_{0}\cdot(K_{1}\cap\cdots\cap K_{s})$. This subgroup contains
$G_{0}$, and therefore corresponds to a subgroup of the quotient $G/G_{0}$,
contained in $(N_{1}/G_{0})\cap(N_{2}/G_{0})\cap\cdots\cap(N_{s}/G_{0})$. This
group can be characterized as the intersection $\sigma(G/N)$ of all non-
trivial normal subgroups in $G/G_{0}$. Since $G_{0}$ is a closed normal
subgroup of the pro-$l$ group $G$, we know that $G/G_{0}$ is an $l$-group, and
hence its center $Z(G/G_{0})$ is non-trivial. Therefore,
$\sigma(G/G_{0})\subseteq Z(G/G_{0})$ is a central subgroup, which implies
that every subgroup of $\sigma(G/G_{0})$ is normal in $G/G_{0}$. By the
definition of $\sigma(G/G_{0})$, this means that $\sigma(G/G_{0})$ has no non-
trivial proper subgroups, and the only such $l$-groups are $\\{e\\}$ and
$\mathbb{Z}/l\mathbb{Z}$. Consequently, either $G_{0}\cdot(K_{1}\cap\cdots\cap
K_{s})/G_{0}$ is trivial, in which case $K_{1}\cap\cdots\cap K_{s}\subseteq
G_{0}$, and we are done, or it is cyclic of order $l$. In this case, we are
given a homomorphism from $\overline{G}_{0}=G_{0}\cdot(K_{1}\cap\cdots\cap
K_{s})$ to $\mathbb{Z}/l\mathbb{Z}$, and therefore a character $\chi_{1}$ of
$\overline{G}_{0}$. Since the group $G$ is totally torsion free, we find that
the group $\overline{G}_{0}^{ab}$ is torsion free, and therefore that the
character group is divisible. Construct a sequence of characters $\chi_{i}$ so
that $\chi_{i}^{l}=\chi_{i-1}$, with $\chi_{1}$ the previously defined
character. This sequence constructs a one-dimensional $l$-BAR of
$\overline{G}_{0}$, which we now induce up to an $l$-BAR ${\eta}$ of dimension
equal to the index of $G_{0}$ in $G$. We claim that the $l$-BAR
$\eta\times\eta_{1}\times\cdots\times\eta_{s}$
is $Ker(\eta)\cap(K_{1}\cap K_{2}\cap\cdots\cap K_{s})$-principal. To see this
we note that the stabilizer of any point in $X_{\eta_{1}}\times\cdots\times
X_{\eta_{s}}$ is equal to $K_{1}\cap\cdots\cap
K_{s}\subseteq\overline{G}_{0}$, and therefore by Proposition 4.8 it suffices
to prove that the restriction of the $G$-action on $X_{\eta}$ to
$K_{1}\cap\cdots\cap K_{s}$ is $Ker(\eta)\cap K_{1}\cap\cdots\cap
K_{s}$-principal. Proposition 4.9 asserts that the restriction to $G_{0}$ of
$\eta$ is $Ker(\eta)$-principal. The result now follows from Proposition 4.8.
$\square$
We need a second fact concerning representations of a profinite $l$-group.
###### Proposition 5.1
Let $\sigma$ denote any irreducible continuous representation of a totally
torsion free profinite $l$-group $G$, over an algebraically closed field of
characteristic prime to $l$. Then $\sigma$ is isomorphic to ${\cal
T}^{s}_{\eta}$ for some $l$-$BAR$ $\eta$ and some $s$.
Proof: The continuity of the $\sigma$ gives that $\sigma$ factors through a
finite quotient $G/N$ of $G$, where $N$ is a closed normal subgroup of finite
index in $G$. By Blichfeldt’s theorem, $\sigma$ (regarded as a representation
of $G/N$) is monomial, and is therefore of the form $i_{K}^{G}(\chi)$ for some
closed subgroup of finite index $K\subseteq G/N$, and some one-dimensional
character $\chi$ of $K$. Let $\overline{K}$ denote the inverse image of $K$ in
$G$. The group $K$ is also totally torsion free, and therefore any
homomorphism from $K$ to a finite cyclic $l$-group can be lifted to a
continuous homomorphism $\hat{\chi}$ from $K$ to $\mathbb{Z}_{l}$. This
homomorphism is a one-dimensional $l$-BAR, and it is clear from the
construction that $\chi$ is equal to $\rho_{s}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}\hat{\chi}$. It now readily
follows that the induced (from $\overline{K}$ to $G$) $l$-BAR is the required
$\eta$. $\square$
We now construct the continuous affine $G$-scheme ${\cal E}G$.
###### Definition 5.1
Let $G$ be any profinite group, which is acting continuously on an affine
scheme $X=Spec(A)$. We say $X$ is cofinite if there is a closed normal
subgroup $N\subseteq G$ of finite index so that the $G$ action on $X$ factors
through the quotient $G/N$. For any cofinite affine $G$-scheme $X$, we define
the subgroup ${\cal N}(X)\subseteq G$ to be the maximal closed normal subgroup
so that the $G$-action on $X$ factors through $G/{\cal N}(X)$. It is clear
that this defines ${\cal N}(X)$ uniquely, and moreover that if we have a map
$X\rightarrow Y$ of cofinite affine $G$-schemes, then ${\cal
N}(X)\subseteq{\cal N}(Y)$. We say a cofinite affine $G$-scheme $X$ is etale
if the group action of $G/{\cal N}(X)$ on $X$ is a free action in the sense of
Proposition 2.6. We define $\mathfrak{C}^{cof}(G)$ to be the category whose
objects are coherent cofinite affine $G$-schemes, and whose morphisms are the
faithfully flat equivariant maps of such $G$-schemes. The full subcategory of
etale actions will be denoted by $\mathfrak{C}^{cof}_{et}(G)$.
###### Proposition 5.2
Suppose that we have a $G$-equivariant map $f:X\rightarrow Y$ of cofinite
etale affine $G$-schemes, which is faithfully flat. Then for any closed
subgroup $K$ of finite index in $G$, the natural map $X/K\rightarrow Y/K$ is
also faithfully flat.
Proof: Since the etale property of an action clearly descends to closed
subgroups, it suffices to prove this for the case $K=G$. We now have the
diagram
$\begin{diagram}$
The morphism $\pi_{Y}\mbox{\hskip 1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$
}}f$ is faithfully flat, since $f$ is and $\pi_{Y}$ is since it is etale and
surjective. $\pi_{X}$ is faithfully is also faithfully flat by the same
reasoning. It follows that $f/G$ is faithfully flat. $\square$
For each closed normal subgroup of finite index $N\lhd G$ we apply Lemma 5.1
to construct an $l$-BAR $\eta_{N}$ which is $N^{\prime}$-principal for some
closed normal subgroup $N^{\prime}\subseteq N$. Let $\mathfrak{N}^{G}$ denote
the partially ordered set of closed normal subgroups of $G$. For each
$(s,N)\in\mathbb{N}\times\mathfrak{N}^{G}$, we let $\theta^{\prime}(s,N)$
denote the (finite) product
$\theta^{\prime}(s,N)=\prod_{\overline{N}\supseteq
N,\overline{N}\in\mathfrak{N}}{\cal T}_{\eta_{\overline{N}}}^{s}$
The construction $\theta^{\prime}$ clearly gives a functor from the category
$\mathbb{N}\times\mathfrak{N}^{G}$ to $\mathfrak{C}^{cof}(G)$, where the
effect on morphisms consists on bonding maps ${\cal
T}_{\eta}^{s+1}\rightarrow{\cal T}_{\eta}^{s}$ together with product
projections. We will need to build a more complicated pro-object. Let
$\mathfrak{R}^{G}$ denote the set of all isomorphism classes of finite
dimensional continuous $k$-linear representations of $G$. For each
$N\in\mathfrak{N}^{G}$, we let $\mathfrak{R}^{G}(N)\subseteq\mathfrak{R}^{G}$
denote the subset of all representations which factor through the quotient
$G/N$. For each $\alpha\in\mathfrak{R}^{G}$, we select an $l$-BAR
$\eta_{\alpha}$ so that for some $s$, $\rho_{s}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}\eta_{\alpha}$ belongs to the
isomorphism class $\alpha$. This is possible because of Proposition 5.1. We
define
$\theta^{\prime\prime}(s,N)=\prod_{\alpha\in\mathfrak{R}^{G}(N)}{\cal
T}_{\eta_{\alpha}}^{s}$
$\theta^{\prime\prime}$ defines a functor from
$\mathbb{N}\times\mathfrak{N}^{G}$ to $\mathfrak{C}^{cof}(G)$, again with
structure maps arising from those in the pro objects ${\cal T}_{\eta}$ and
product projections. We let ${\cal E}G$ denote the inverse limit of
$\theta=\theta^{\prime}\times\theta^{\prime\prime}$ in the category of affine
$G$-schemes. ${\cal E}G$ is thus an affine scheme whose affine coordinate ring
is the colimit of the rings $A(\theta(s,N))$ over the category
$\mathbb{N}\times\mathfrak{N}^{G}$, where $A(X)$ denotes the affine coordinate
ring of $X$.
###### Proposition 5.3
For any pair $N\subseteq N^{\prime}\subseteq G$ of closed normal subgroups of
finite index, and any $s\leq s^{\prime}$, the map of schemes
$\theta(s^{\prime},N^{\prime})\rightarrow\theta(s,N)$ is faithfully flat. It
follows that ${\cal E}G$ is coherent.
Proof: On affine coordinate rings, each map is of the form
$A=k[x_{1}^{\pm\frac{1}{l^{s}}},\ldots,x_{t}^{\pm\frac{1}{l^{s}}}]\hookrightarrow
k[x_{1}^{\pm\frac{1}{l^{s^{\prime}}}},\ldots,x_{t^{\prime}}^{\pm\frac{1}{l^{s^{\prime}}}}]=B$
where $t^{\prime}\geq t$. In this case, it is an easy computation that $B$ is
a free module over $A$. Both rings are Noetherian, therefore coherent, and the
result now follows from Proposition 2.3. $\square$
###### Proposition 5.4
The $G$-action on the affine coordinate ring $A[{\cal E}G]$ is continuous,
when $A[{\cal E}G]$ is equipped with the discrete topology.
Proof: Each ring $A[\theta(s,N)]$ is acted on by $G$ via an action which
factors through a finite quotient of $G$. Since $A[{\cal E}G]$ is the colimit
of the rings $A[\theta(s,N)]$, the result follows. $\square$
We will now prove that the $G$-action on ${\cal E}G$ is etale. Let
$\overline{\theta}$ denote the functor $(s,N)\rightarrow\theta(s,N)/N$.
$\overline{\theta}(s,N)$ is equipped with a $G/N$-action.
###### Lemma 5.2
The natural map
$\mathrel{\mathop{lim}\limits_{\leftarrow}}\theta\rightarrow\mathrel{\mathop{lim}\limits_{\leftarrow}}\overline{\theta}$
is an isomorphism of affine schemes.
Proof: We prove that the map induces an isomorphism on affine coordinate
rings. The map
$A[\mathrel{\mathop{lim}\limits_{\leftarrow}}\overline{\theta}]\rightarrow
A[\mathrel{\mathop{lim}\limits_{\leftarrow}}\theta]$ is clearly injective,
since for every $(s,N)$, the map $A[\overline{\theta}(s,N)]\rightarrow
A[\theta(s,N)]$ can be identified with the inclusion
$A[\theta(s,N)]^{N}\rightarrow A[\theta(s,N)]$. Given an element $x\in
A[\mathrel{\mathop{lim}\limits_{\leftarrow}}\theta]$, it is in the image of
$A[\theta(s,N)]$ for some $s$ and $N$. This means that it is fixed by a normal
subgroup of finite index $N^{\prime}\subseteq N$, and consequently that it is
in the image of
$A[\theta(s,N^{\prime})]^{N^{\prime}}=A[\overline{\theta}(s,N^{\prime})]$.
This proves surjectivity. $\square$
###### Lemma 5.3
Let $N\subseteq G$ be a closed subgroup of finite index. Then
$\mathrel{\mathop{lim}\limits_{\mathrel{\mathop{\leftarrow}\limits_{s}}}}\overline{\theta}(s,N)$
an etale $G/N$-action for all $N$.
Proof: Proposition 4.6 shows that
$\mathrel{\mathop{lim}\limits_{\mathrel{\mathop{\leftarrow}\limits_{s}}}}{\theta}(s,N)$
is an etale action of $G/N^{\prime}$ for some closed normal subgroup
$N^{\prime}\subseteq N$. It follows that the action of $G/N$ on
$\mathrel{\mathop{lim}\limits_{\mathrel{\mathop{\leftarrow}\limits_{s}}}}{\theta}(s,N)/N=\mathrel{\mathop{lim}\limits_{\mathrel{\mathop{\leftarrow}\limits_{s}}}}\overline{\theta}(s,N)$
is etale. $\square$
###### Corollary 5.1
For every closed normal subgroup $N\subseteq G$ of finite index, the
$G/N$-action on $\overline{\theta}(s^{\prime},N)$ is etale for sufficiently
large $s$.
Proof: We let $N^{\prime}\subseteq N\subseteq G$ be the closed normal subgroup
of $G$ so that $\eta_{N}$ is $N^{\prime}$-principal, and therefore that
$\mathrel{\mathop{lim}\limits_{\mathrel{\mathop{\leftarrow}\limits_{s}}}}{\theta}(s,N)$
is an etale $G/N^{\prime}$-action. It follows from Proposition 4.7 that there
is an $s\geq 0$ so that the action of
$G/\eta^{-1}(l^{s^{\prime}}\mathfrak{T})$ on
$\mathrel{\mathop{lim}\limits_{\mathrel{\mathop{\leftarrow}\limits_{s}}}}{\theta}(s,N)/(\eta^{-1}(l^{s^{\prime}}\mathfrak{T}))$
is etale for all $s^{\prime}\geq s$. Because $N$ is of finite index, there is
also a $t\geq 0$ so that $\eta^{-1}(l^{t^{\prime}}\mathfrak{T})\subseteq N$
for all $t^{\prime}\geq t$. Letting $S=max(s,t)$, we have that for every
$s^{\prime}\geq S$, the $G/N$-action on $\overline{\theta}(s^{\prime},N)$ is
defined as the natural action of $G/N$ on the
$N/\eta^{-1}(l^{s^{\prime}}\mathfrak{T})$-orbits of the action of the
$G/\eta^{-1}(l^{s^{\prime}}\mathfrak{T})$-action on
$\mathrel{\mathop{lim}\limits_{\mathrel{\mathop{\leftarrow}\limits_{s}}}}{\theta}(s,N)/(\eta^{-1}(l^{s^{\prime}}\mathfrak{T}))$.
Since the $G/\eta^{-1}(l^{s^{\prime}}\mathfrak{T})$-action is etale, so is the
$G/N$-action. $\square$
###### Proposition 5.5
${\cal E}G$ is etale.
Proof: We must show that for any closed subgroup of finite index $N$, the
$G/N$ action on $({\cal E}G)/N$ is etale. We have that $({\cal E}G)/N$ is the
inverse limit of the functor $(s,N)\rightarrow\theta(s,N)/N$. This is a
diagram, parametrized by a directed set, with values in the category of affine
schemes with $G/N$-action. For every closed normal subgroup of finite index
$N^{\prime}\subseteq G$, let $s(N^{\prime})\geq 0$ be an integer so that the
$G/N$-action on $\theta(s,N^{\prime})/N$ is etale for all $s\geq
s(N^{\prime})$. The subset
$\hat{\mathfrak{N}}=\\{(s,N^{\prime})\mbox{ such that }s\geq s(N^{\prime})\\}$
contained in $\mathbb{N}\times\mathfrak{N}^{G}$ is clearly final, and so
${\cal E}G/N$ can be described as the inverse limit of the restriction of the
functor $N^{\prime}\rightarrow\theta(s,N^{\prime})/N$ to the subset
$\hat{\mathfrak{N}}$. Proposition 2.8 now gives the desired result. $\square$
###### Proposition 5.6
Let $W$ denote any continuous affine $G$-scheme, and assume that $A[W]$ is
Noetherian approximable. Then ${\cal E}G\times W$ is coherent, and the
$G$-action is etale. The natural map
$K({\cal E}G\mathrel{\mathop{\times}\limits_{G}}W)\rightarrow K^{G}({\cal
E}G\times W)$
therefore exists and is an equivalence.
Proof: ${\cal E}G\times W$ clearly is an etale action, using the base change
property for etale actions. The affine coordinate ring of ${\cal E}G\times W$
can be written as the colimit of rings of the form
$A[W][x_{1}^{\pm\frac{1}{l^{s}}},\ldots,x_{t}^{\pm\frac{1}{l^{s}}}]$, which
are coherent by Proposition 2.4 since $A[W]$ is Noetherian approximable.
Further, each of the inclusions in the colimit system is of the form
$A[W][x_{1}^{\pm\frac{1}{l^{s}}},\ldots,x_{t}^{\pm\frac{1}{l^{s}}}]\hookrightarrow
A[W][x_{1}^{\pm\frac{1}{l^{s^{\prime}}}},\ldots,x_{t^{\prime}}^{\pm\frac{1}{l^{s^{\prime}}}}]$
and are therefore clearly faithfully flat. This demonstrates the coherence of
$A[{\cal E}G\times W]$. The result now follows by Proposition 2.16. $\square$
Remark: The conditions on $\theta^{\prime}$ provide a freeness property of
${\cal E}G$. Extending by creating the product with $\theta^{\prime\prime}$
enlarges the pro-scheme so as to trivialize the action of $R[G]$ on the
equivariant $K$-groups of ${\cal E}G$.
We prove the above mentioned technical result on the action of $R[G]$ on
$K^{G}_{*}({\cal E}G)$. which will be useful in analyzing completions of
$K^{G}({\cal E}G)$.
###### Proposition 5.7
Let $G$ be a totally torsion free $l$-profinite group. Then the action of
$\pi_{0}K^{G}(k)\cong R[G]$ satisfies $I_{G}\cdot\pi_{*}K^{G}({\cal
E}G)=\\{0\\}$.
Proof: It clearly suffices show that $j(I_{G})=\\{0\\}$. This means that we
will need to show that for any representation $\rho$ of $G$, we have that
$[\rho]-dim(\rho)$ vanishes in $K^{G}_{0}({\cal E}G)$. Blichfeldt’s theorem
(see [35]) tells us that any irreducible representation of $G$ is monomial,
i.e it is isomorphic to an induced representation $i_{G_{0}}^{G}(\chi)$ where
$G_{0}$ is a subgroup of index $dim(\rho)$ and $\chi$ is a one dimensional
character of $G_{0}$. Because of the totally torsion free property of $G$,
$\chi$ is of the form $\rho_{s}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}\eta_{\chi}$ where
$\eta_{\chi}$ is a one-dimensional $l$-BAR of $G_{0}$, and the $\rho_{s}$’s
are defined in Section 4. Inducing $\eta_{\chi}$ up to $G$ we obtain an
$l$-BAR $\eta$ of $G$, and it is clear from the definition that $\rho$ can be
identified with $\rho_{s}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}\eta$. The conclusion is that
every irreducible representation of $G$ is of the form $\rho_{s}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}\eta$ for some $l$-BAR of $G$,
and some $s$. Next, we observe that any $l$-BAR is a factor of a principal
one. This is an immediate consequence of Lemma 5.1. For any $l$-BAR $\eta$ we
have the associated permutation representation $\overline{\eta}$. We will
first prove that $[\rho_{s}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}\eta]-[\rho_{s}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}\overline{\eta}]=0$ in
$K_{0}^{G}({\cal E}G)$. Fix $s\geq 0$. Given an $l$-BAR $\eta$, we let
$K_{\eta}$ denote the kernel of $\rho_{s}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}\eta$. Similarly, let
$\eta^{\prime}$ be a $K_{\eta^{\prime}}$-principal $l$-BAR containing $\eta$
as a factor, where $K_{\eta^{\prime}}\subseteq K_{\eta}$. We now have a
factorization
$R[G/K_{\eta}]\cong K_{0}^{G/K_{\eta}}(k)\rightarrow K_{0}^{G/K_{\eta}}({\cal
T}^{s}_{\eta})\stackrel{{\scriptstyle
j}}{{\longrightarrow}}K_{0}^{G/K_{\eta^{\prime}}}({\cal
T}^{s}_{\eta^{\prime}})\rightarrow K_{0}^{G}({\cal E}G)$
since ${\cal T}_{\eta}^{s}$ is one of the factors occurring in the inverse
system defining ${\cal T}_{\eta}$, and therefore one of the factors occurring
in ${\cal E}G$. The map $j$ is the composite
$K_{0}^{G/K_{\eta}}{\cal T}_{\eta}^{s}\rightarrow
K_{0}^{G/K_{\eta^{\prime}}}{\cal T}_{\eta}^{s}\rightarrow
K_{0}^{G/K_{\eta^{\prime}}}{\cal T}_{\eta^{\prime}}^{s}$
where the left hand map is the inflation map associated to the projection
$G/K_{\eta^{\prime}}\rightarrow G/K_{\eta}$, and the right hand map is the map
induced by the equivariant projection ${\cal
T}_{\eta^{\prime}}^{s}\rightarrow{\cal T}_{\eta}^{s}$. Letting
$\rho=\rho_{s}\mbox{\hskip 1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$
}}\eta$ and $\overline{\rho}=\rho_{s}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}\overline{\eta}$ it will
suffice to prove that the image of $[\rho]-[\overline{\rho}]=0$ in
$K_{0}^{G}\mathbb{T}(k,\rho)$ vanishes. Let $N$ be the kernel of $\rho$, a
normal subgroup of finite index in $G$. Note that $\rho$ and $\overline{\rho}$
may be regarded as elements in $K^{G/N}_{0}\mathbb{A}(k,\rho)$. We observe
that there is a factorization
$R[G/N]\cong K^{G/N}(k)\rightarrow K^{G/N}_{0}\mathbb{A}(k,\rho)\rightarrow
K_{0}^{G}\mathbb{T}(k,\rho)\mbox{ }(=K_{0}^{G}{\cal T}_{\eta}^{s})$
We now select a basis $\\{x_{1},\ldots,x_{n}\\}$ for the representation space
of $\rho$, where $n=dim(\rho)$, so that $\rho$ is monomial with respect to
this basis, and we let $\overline{\rho}$ denote the associated permutation
representation. The group $K^{G/N}_{0}\mathbb{A}(k,\rho)$ is the group
completion of the commutative monoid of isomorphism classes of finitely
generated left modules over the skew polynomial ring
$A_{\rho}=k[x_{1},\ldots,x_{n}]_{\rho}[G/N]$, by Proposition 2.17. We
construct a short exact sequence of left $A_{\rho}$-modules as follows. First,
let $G_{0}\subseteq G$ be the stabilizer of the line spanned by the element
$x_{1}$, so $\rho$ is induced from a one-dimensional representation of
$G_{0}/N$, so that $\rho$ is induced from a one-dimensional representation
$\chi$ of $G_{0}$. We let $F$ be the module
$k[x_{1}]_{\chi}[G_{0}/N]\mathrel{\mathop{\otimes}\limits_{k[G_{0}/N]}}\epsilon$,
where $\epsilon$ denotes the trivial representation of $G_{0}/N$. The module
$F$ is free of rank one over $k[x_{1}]$, and we may write $F=k[x_{1}]b$, where
$b=1\otimes 1$ is the basis element which is acted on trivially by $G_{0}/N$.
We let $k$ denote the left $k[x_{1}]_{\chi}[G_{0}/N]$-module which as a
$k$-module is cyclic of rank one, and on which all elements of $G_{0}/N$ act
by the identity and such that $x_{1}\cdot k=0$. There is a natural
homomorphism $\pi:F\rightarrow k$, which is the identification $F/xF\cong k$.
Let $F^{\prime}$ denote the kernel of $\pi$. It is the free cyclic
$k[x]$-module generated by $xb$, and from this description it is clear that
$F^{\prime}$ is isomorphic to
$k[x_{1}]_{\chi}[G_{0}/N]\mathrel{\mathop{\otimes}\limits_{k[G_{0}/N]}}\chi$.
The conclusion is that the equation $[k]=[\epsilon]-[\chi]$ holds in
$K_{0}^{G_{0}/N}\mathbb{A}(k,\chi)\cong R[G_{0}/N]$, where we are interpreting
this group as the projective class group of finitely generated modules over
$k[x_{1}]_{\chi}[G_{0}/N]$. Performing induction from $G_{0}/N$, and using
Proposition 2.18, we find that the equation
$[M]=[i_{G_{0}/N}^{G/N}(\epsilon)]-[i_{G_{0}/N}^{G/N}(\chi)]$
holds in $K_{0}^{G/N}\mathbb{A}(k,\rho)\cong R[G/N]$, where
$M=k[x_{1},\ldots,x_{n}]_{\rho}[G/N]\mathrel{\mathop{\otimes}\limits_{K[x_{1}]_{\chi}[G_{0}/N]}}k$
We observe that $i_{G_{0}/N}^{G/N}(\epsilon)$ is the permutation
representation $\overline{\rho}$ of $G$. It is also readily checked that the
element $x_{1}x_{2},\cdots x_{n}$ annihilates the module $M$, and therefore
that it becomes trivial on tensoring up to $k[x_{1}^{\pm 1},\ldots,x_{n}^{\pm
1}]$, which means that the equation
$0=[i_{G_{0}/N}^{G/N}(\epsilon)]-[i_{G_{0}/N}^{G/N}(\rho]$
holds in $K_{0}^{G/N}\mathbb{T}(k,\rho)\cong R[G/N]$. This gives the fact that
$[\rho]-[\overline{\rho}]=0$ holds in $K_{0}({\cal E}G)$ for every irreducible
representation $\rho$ of $G$. Since every permutation representation of a
finite group is decomposable (it contains a copy of the trivial representation
of $G$), this means that for every irreducible representation $\rho$ of $G$,
the class $[\rho]\in K_{0}({\cal E}G)$ is equal to a sum of images of
representations of smaller dimension. Note that the special case $n=1$ of this
result shows that one-dimensional representations $\chi$ satisfy the condition
that $[\chi]-dim(\chi)$ vanishes in $K^{G}({\cal E}G)$. An induction on $n$
now shows that $[\rho]$ is equal to $dim(\rho)$, which was to be shown.
$\square$
###### Corollary 5.2
Let $W$ be any affine scheme with continuous $G$-action. Then the action of
$R[G]$ satisfies $I_{G}\cdot K^{G}({\cal E}G\times W)=\\{0\\}$.
Proof: This follows since $K^{G}({\cal E}G\times W)=\\{0\\}$ is a ring
spectrum, and we have that $I_{G}$ maps trivially into $\pi_{0}K^{G}({\cal
E}G\times W)$ since there is a factorization
$R[G]\rightarrow K^{G}({\cal E}G)\rightarrow K^{G}({\cal E}G\times W)$
and we already know from Proposition 5.7 that the map $R[G]\rightarrow
K^{G}({\cal E}G)$ is trivial. $\square$
Now we consider the diagram of commutative ring spectra
$\begin{diagram}$
where $r$ denotes the ring map induced by the fundamental class in the sphere
spectrum, and $\epsilon$ is the augmentation. The functorality of the derived
completion construction means that the diagram induces a map
$\beta:K^{G}({\cal E}G\times W)^{\wedge}_{l}=K^{G}({\cal E}G\times
W)^{\wedge}_{r}\rightarrow K^{G}({\cal E}G\times W)^{\wedge}_{\epsilon}$
###### Corollary 5.3
The map $\beta$ is an equivalence of spectra.
Proof: The algebraic-to-geometric spectral sequence (Theorem 7.1 of [7]) and
the preceding Proposition show that it suffices to prove that the map on
derived completions $\hat{\beta}:M^{\wedge}_{l}\rightarrow
M^{\wedge}_{\epsilon}$ induced by the diagram
$\begin{diagram}$
is an equivalence for modules $M$ with trivial action by $I_{G}\subseteq
R[G]$, or equivalently which are obtained by restricting scalars along the
homomorphism $\epsilon:R[G]\rightarrow\mathbb{Z}$. In order to prove this, we
regard the ring $R[G]$ as a commutative ring spectrum, and let
$\tilde{\mathbb{Z}}$ denote a cofibrant replacement for the $R[G]$-algebra
$\mathbb{Z}$, with algebra structure given by $\epsilon$, in the model
category of commutative ring spectra discussed in [36]. We have the natural
structure map $R[G]\rightarrow\tilde{\mathbb{Z}}$, as well as a homomorphism
$\tilde{\mathbb{Z}}\rightarrow R[G]$ of commutative ring spectra obtained by
extending the homomorphism of rings $\mathbb{Z}\rightarrow R[G]$ to the
cofibrant replacement $\tilde{\mathbb{Z}}$. We also have a homomorphism
$\tilde{\mathbb{Z}}\rightarrow\mathbb{H}_{l}$ which is compatible with the
mod-$l$ reduction of the augmentation on $R[G]$. We now have two monads $S$
and $T$ on the category of $\tilde{\mathbb{Z}}$-module spectra, where
$S(M)=\mathbb{H}_{l}\mathrel{\mathop{\wedge}\limits_{\tilde{\mathbb{Z}}}}M$
and $T(M)=\mathbb{H}_{l}\mathrel{\mathop{\wedge}\limits_{R[G]}}M$. Further,
there are natural transformations $T\rightarrow S$ and $S\rightarrow T$
induced by the ring homomorphisms $\tilde{\mathbb{Z}}\rightarrow
R[G]\rightarrow\tilde{\mathbb{Z}}$. It follows that the natural
transformations of triples induce equivalence on the total spectra of the
cosimplicial objects of the triples $S$ and $T$, by Theorem 2.15 of [7]. These
total spectra attached to $S$ and $T$ are the derived completions over
$\tilde{\mathbb{Z}}$ and $R[G]$ respectively. The fact that weak equivalences
of commutative ring spectra induce equivalences of derived completions now
gives the result. $\square$
We also need a similar result for $K^{G}(\overline{F})$, where $G$ is the
absolute Galois group of a field $F$, containing $k$ as a subfield.
###### Proposition 5.8
With $G,F,\overline{F}$ as above, the action of $R[G]$ on
$K^{G}_{*}(\overline{F})$ is trivial in the sense that $I_{G}\cdot
K^{G}_{*}(\overline{F})=\\{0\\}$.
Proof: It is proved in [22] that the category of finitely generated
$G$-twisted $\overline{F}$-modules is equivalent to the category of finite
dimensional vector spaces over $F$, and that the equivalence can be taken by
applying the functor
$V\rightarrow\overline{F}\mathrel{\mathop{\otimes}\limits_{F}}V$. There is a
natural twisted $G$-action on
$\overline{F}\mathrel{\mathop{\otimes}\limits_{F}}V$. Since $\pi_{0}K(F)$ is a
free group of rank one, given only by the dimension, it is therefore clear
that for any finite dimensional $k$-linear representation $(W,\rho)$ of $G$,
$\overline{F}\mathrel{\mathop{\otimes}\limits_{F}}W$ is isomorphic to
$\overline{F}\mathrel{\mathop{\otimes}\limits_{F}}W^{0}$, where $W^{0}$
denotes $W$ equipped with trivial action. This gives the result. $\square$
###### Corollary 5.4
There is a natural equivalence $K^{G}(\overline{F})^{\wedge}_{\epsilon}\cong
K(F)^{\wedge}_{l}$.
Proof: The proof is identical to that given in Proposition 5.3 above. That
proof uses only the triviality of the $I_{G}$-action on $\pi_{0}$, and so
applies here as well. $\square$
We conclude this section by proving that the equivariant $K$-theory of ${\cal
E}G\times W$ can be computed as a particular kind of colimit of equivariant
$K$-theory spectra of smooth schemes of finite type over $k$. This will be
important to us because the rigidity results we prove in [10] will apply only
to such systems.
###### Definition 5.2
Let $G$ be a profinite group, and as usual let $\mathfrak{N}^{G}$ denote the
directed set of closed normal subgroups of finite index in $G$. Let
$F:\mathfrak{N}^{G}\rightarrow\mathfrak{C}^{cof}_{et}(G)$ denote any functor
($\mathfrak{C}^{cof}_{et}(G)$ was defined in Definition 5.1 above). We say $F$
is adapted if ${\cal N}(F(N))=N$ for all closed normal $N\subseteq G$. Let
$\rho$ be any continuous homomorphism from $G$ to the absolute Galois group of
a field $K$. Let $L$ denote the absolute closure of $K$. We define spectrum
valued functors $K(F,L)$ and $K^{G}(F,L)$ associated to this data by
$K(F,L)=\mathrel{\mathop{hocolim}\limits_{N\in\mathfrak{N}^{G}}}K(F(N)\mathrel{\mathop{\times}\limits_{G/N}}L^{N})$
and
$K^{G}(F,L)=\mathrel{\mathop{hocolim}\limits_{N\in\mathfrak{N}^{G}}}K^{G/N}(F(N)\times
L^{N})$
These colimits are naturally equivalent by Proposition 2.6. $F$ is said to be
smooth (respectively of finite type) if each $F(N)$ is smooth (respectively of
finite type).
We are now able to express the $K$-theory of ${\cal E}G$ in these terms.
###### Proposition 5.9
Let $G$ be a totally torsion free profinite $l$-group, and let $\rho$ be a
homomorphism into the absolute Galois group of a field $K$. Let $L$ denote the
algebraic closure of $K$. Then there is a smooth adapted functor
$F(G,\rho):\mathfrak{N}^{G}\rightarrow\mathfrak{C}^{cof}_{et}(G)$ of finite
type, and an equivalence $K^{G}(F(G,\rho),L)\rightarrow K^{G}({\cal E}G\times
Spec(L))$. This equivalence is natural for inclusions of fields with
continuous $G$-action.
Proof: For every closed normal subgroup of finite index $N\subseteq G$,
Corollary 5.1 asserts that there is an integer $s(N)$ so that the action of
$G/N$ on $\overline{\theta}(s^{\prime},N)$ for every $s\geq s(N)$, is etale.
We now define a function
$j:\mathfrak{N}^{G}\rightarrow\mathbb{N}\times\mathfrak{N}^{G}$ by
$j(N)=\mathrel{\mathop{max}\limits_{N^{\prime}\subseteq N}}s(N)$
It is clear from the definition that
$J:\mathfrak{N}^{G}\rightarrow\mathbb{N}\times\mathfrak{N}^{G}$, defined by
$J(N)=(j(N),N)$ is a map of directed sets, that $\overline{\theta}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}J$ is an adapted functor, and
further that the image of $J$ is a final subset of $\mathfrak{N}^{G}$. It
follows that we have the functors $K^{G}(\overline{\theta}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}J,L)$ for any $\rho$ and $L$ as
in the statement of the Proposition. Moreover, we have an equivalence of
colimits from $K^{G}(\overline{\theta}\mbox{\hskip
1.42271pt}\raisebox{1.13791pt}{\tiny{$\circ$ }}J,L)$ to $K^{G}({\cal E}G\times
Spec(L))$. Finally, each of the schemes $F(N)$ are of finite type and smooth,
since they are quotients by a finite etale group action on a torus. $\square$
Remark: Note that the cases we are interested in are $K=L=k$ and $K=F$,
$L=\overline{F}$.
We write $F^{G}_{L}$ for $l$-adic completion of the the functor
$N\rightarrow\theta(\nu(N))$ given in the proof of the Proposition above. We
will prove that the maps $K^{G}(k)\rightarrow K^{G}({\cal E}G)$ and
$K^{G}(Spec(\overline{F}))\rightarrow K^{G}({\cal E}G\times
Spec(\overline{F}))$ induce equivalences on derived completions along the
homomorphism $\epsilon:K^{G}(k)\rightarrow\mathbb{H}_{l}$, and further that
the completions of $K^{G}({\cal E}G)$ and $K^{G}({\cal E}G\times
Spec(\overline{F}))$ are just $l$-adic completions. Consider the diagram
$\begin{diagram}$
We will be proving that the left hand horizontal arrows are equivalences. The
right hand horizontal arrows are equivalences by Proposition 5.9 above. It
follows that in order to prove that the left hand vertical arrow is an
equivalence, it will suffice to prove that the right hand vertical arrow is
one. This is what will be proved in [10].
## 6 Algebraic properties of certain representation rings
In this section, we analyze some particular properties of the representation
rings of certain semidirect products of symmetric groups with profinite groups
of the form $\mathbb{Z}_{l}^{n}$. The results of this section are the key
technical tools which allow us to prove our main results.
###### Theorem 6.1
Let $k$ be a field of characteristic $l$, and let $\epsilon:A\rightarrow k$ be
an augmented $k$-algebra with augmentation ideal $I_{A}$. Let $G\subseteq
G^{\prime}\subseteq\Sigma_{n}$ denote subgroups of $l$-power order. Let
$\displaystyle B=\mathop{\otimes}^{n}_{k}A$, equipped with the action of
$\Sigma_{n}$, and therefore of $G$ and $G^{\prime}$. The algebra $B$ is given
the tensor product augmentation $\epsilon_{B}$, and the algebras $B^{G}$ and
$B^{G^{\prime}}$ both obtain augmentations by restriction of $\epsilon_{B}$,
and we denote the restrictions by $\epsilon_{B^{G}}$ and
$\epsilon_{B^{G^{\prime}}}$ respectively. Further, we let $I_{B^{G}}$ and
$I_{B^{G^{\prime}}}$ denote the corresponding augmentation ideals. Finally,
suppose that every element in $A$ has an $l$-th root. Then the equality of
ideals
$I_{B^{G}}=I_{B^{G^{\prime}}}\cdot B^{G}$
holds.
Proof: We first prove that the $k$-algebra $B^{G}$ also has the property that
every element admits an $l$-th root. To see this, we let $\mathfrak{A}$ denote
a $k$-basis for $A$. Then we form the product $\mathfrak{B}=\mathfrak{A}^{n}$,
note that it forms a $k$-basis for $B$, and let $G$ act on $\mathfrak{B}$ via
the inclusion $G\hookrightarrow\Sigma_{n}$. For every orbit
$\mathfrak{o}=\\{\beta_{1},\ldots,\beta_{s}\\}$, we let $t(\mathfrak{o})$
denote the element $\beta_{1}+\cdots+\beta_{s}\in B^{G}$. It is clear that the
elements $t(\mathfrak{o})$, as $\mathfrak{o}$ ranges over the orbits of the
$G$ action on $\mathfrak{B}$, form a $k$-basis for $B^{G}$, and therefore that
it suffices to prove that each element $t(\mathfrak{o})$ admits an $l$-th root
in $B^{G}$. We suppose that we are given an orbit $\mathfrak{o}$, with orbit
representative $\beta\in\mathfrak{o}$. Let $G_{\beta}$ denote the stabilizer
of $\beta$. The group $G_{\beta}$ now acts on the set $\\{1,\ldots,n\\}$ via
its inclusion into $\Sigma_{n}$, and we let
$\\{1,\ldots,n\\}=S_{1}\cup S_{2}\cup\cdots\cup S_{l}$
denote the orbit decomposition of this action. Since $G_{\beta}$ stabilizes
$\beta$, it is clear that there are elements $\\{\alpha_{j}\\}_{j=1}^{l}$,
with $\alpha\in\mathfrak{A}$, so that
$\beta=\alpha_{j(1)}\otimes\alpha_{j(2)}\otimes\cdots\otimes\alpha_{j(n)}$
where $j(i)$ denotes the integer so that $i\in S_{j(i)}$. For each
$j=1,\ldots,l$, we select an element $\overline{\alpha}_{j}$ so that
$(\overline{\alpha}_{j})^{l}=\alpha_{j}$. Also, select left coset
representatives $\\{g_{1},\ldots,g_{k}\\}$ for the collection of left cosets
$G/G_{\beta}$. The element
$\sum_{m=1}^{k}g_{m}(\overline{\alpha}_{j(1)}\otimes\overline{\alpha}_{j(2)}\otimes\cdots\otimes\overline{\alpha}_{j(l)})$
is now clearly an element in $B^{G}$, and it is equally clearly an $l$-th root
of $t(\mathfrak{o})$, which gives us the claim. We now proceed to the proof of
the theorem.
We consider first the case where $G\lhd G^{\prime}$ and where the index of $G$
in $G^{\prime}$ is $l$. Consider any element $b\in I_{B^{G}}$. We let
$\overline{b}$ denote any $l$-th root of $b$ in $B^{G}$. It is immediate that
$\overline{b}\in I_{B^{G}}$ as well. Now construct the polynomial
$\varphi(X)=\prod_{\gamma\in G^{\prime}/G}(X-\overline{b}^{\gamma})$
Note that the notation $\overline{b}^{\gamma}$ is well defined for $\gamma\in
G^{\prime}/G$ because $\overline{b}\in B^{G}$. From the definition of
$\varphi$, it is clear that $\varphi(\overline{b})=0$, and that $\varphi$ is
of the form
$\varphi(X)=X^{l}+\sum_{i=1}^{l}(-1)^{i}c_{i}X^{l-i}$
where $c_{i}$ is the $i$-th elementary symmetric function in the elements
$\overline{b}^{\gamma}$, as $\gamma$ varies over all the elements of
$G^{\prime}/G$. Note that $c_{i}\in I_{B^{G^{\prime}}}$. Since
$\overline{b}^{l}=b$, it follows that
$b=\sum_{i=1}^{l}(-1)^{i+1}c_{i}\overline{b}^{l-i}$
which gives the result in this case. For the general case, we may form a
sequence of subgroups
$G=G_{t}\subseteq G_{t-1}\subseteq G_{t-2}\subseteq\cdots\subseteq
G_{1}\subseteq G_{0}=G^{\prime}$
where each inclusion is the inclusion of a normal subgroup of index $l$. This
is an immediate consequence of the fact that $G^{\prime}$ is an $l$-group. The
result for the general case can now be obtained from the case of a normal
subgroup of index $l$ via an induction on $t$. We wish to prove that
$I_{B^{G_{0}}}B^{G_{t}}=I_{B^{G_{t}}}$, and by induction we suppose that we
have proved that $I_{B^{G_{0}}}B^{G_{t-1}}=I_{B^{G_{t-1}}}$ and
$I_{B^{G_{t-1}}}B^{G_{t}}=I_{G_{t}}$. We obtain
$I_{B^{G_{t}}}=I_{B^{G_{t-1}}}B^{G_{t}}=I_{B^{G_{0}}}B^{G_{t-1}}B^{G_{t}}=I_{B^{G_{0}}}B^{G_{t}}$
$\square$
We now obtain the result which allows us to apply Corollary 2.6.
###### Theorem 6.2
Let $k$ be a field of characteristic $l$, and $A$ an augmented $k$-algebra in
which every element has an $l$-th root. Let $G\subseteq
G^{\prime}\subseteq\Sigma_{n}$ be an inclusion of $l$-groups, and let
$\displaystyle B=\mathop{\otimes}^{n}_{k}A$ be equipped with the natural
action of $\Sigma_{n}$ and therefore of $G$ and $G^{\prime}$. Also equip $B$
with the tensor product augmentation. Suppose we have a diagram of augmented
(to $k$) rings
$\begin{diagram}$
in which $\pi$ and $\pi^{\prime}$ are surjective, and so that the kernel
$\mathfrak{I}_{R}$ of $\pi$ satisfies $(\mathfrak{I}_{R})^{N}=\\{0\\}$ for
some positive integer $N$. Then there is an integer $M$ so that
$(I_{R})^{M}\subseteq I_{R^{\prime}}\cdot R$, where $I_{R}$ and
$I_{R^{\prime}}$ denote the augmentation ideals of $R$ and $R^{\prime}$
respectively.
Proof: Let $r\in I_{R}$. Then $\pi(r)\in I_{B^{G}}$, and consequently, by
Proposition 6.1, $\pi(r)\in I_{B^{G^{\prime}}}B^{G}$. We can therefore write
$\pi(r)=\sum_{s}i(y_{s})z_{s}$
where $y_{s}\in I_{B^{G^{\prime}}}$ and $z_{s}\in B^{G}$. Since $\pi$ and
$\pi^{\prime}$ are surjective, we can write
$y_{s}=\pi^{\prime}(\overline{y}_{s})$ and $z_{s}=\pi(\overline{z}_{s})$ for
each $s$. We let
$\overline{r}=\sum_{s}f(\overline{y}_{s})\overline{z}_{s}$
It is clear from the construction that $\pi(r-\overline{r})=0$, and of course
that $\overline{r}\in I_{R^{\prime}}R$. We now let $M$ be any power of $l$
which is greater than $N$. Then we have
$0=(r-\overline{r})^{M}=r^{M}-\overline{r}^{M}$
so $r^{M}=\overline{r}^{M}$, which gives that $r^{M}\in I_{R^{\prime}}R$.
$\square$
We will use this result to study certain representation rings. A first step is
to study the behavior of the mod-$l$ representation rings of certain wreath
products. Recall that for any groups $G$ and $H$, with $G$ finite, the wreath
product $G\wr H$ is the semidirect product $G\ltimes H^{\\#(G)}$, where the
$G$ action on $H^{\\#(G)}$ is by permutation of coordinates. For any group
$G$, we will write $\mathfrak{R}_{l}[G]$ for
$R[G]\mathrel{\mathop{\otimes}\limits_{\mathbb{Z}}}\mathbb{F}_{l}$.
###### Proposition 6.1
Let $C_{l}$ denote the cyclic group of order $l$, where $l$ is a prime. Then
for any profinite group $K$, we construct $C_{l}\wr K$, with the inclusion
$i:K^{l}\hookrightarrow C_{l}\wr K$. The then following results hold.
1. 1.
The restriction map $i^{*}:\mathfrak{R}_{l}[C_{l}\wr
K]\rightarrow\mathfrak{R}_{l}[K^{l}]$ is surjective onto the invariant subring
$\mathfrak{R}_{l}[K^{l}]^{C_{l}}$.
2. 2.
The kernel $\mathfrak{I}$ of $i^{*}$ satisfies $\mathfrak{I}^{l}=0$.
Proof: Consider first the case where $K$ is finite. Then the results of
section 4.3 of [24] show that an additive basis of $\mathfrak{R}_{l}[C_{l}\wr
K]$ can be constructed as follows. Let $\mbox{Irr}(K)$ denote the set of
irreducible representations of $K$. Then we have the permutation action of
$C_{l}$ on $\mbox{Irr}(K)^{l}$, which can then be decomposed into orbits,
which either consist of one element or of $l$ elements. Letting the orbit set
be denoted by $\mathfrak{O}$, we write
$\mathfrak{O}=\mathfrak{O}_{(1)}\coprod\mathfrak{O}_{(l)}$ for this
decomposition. A basis for $\mathfrak{R}_{l}[C_{l}\wr K]$ is now in bijective
correspondence with the set
$\mathfrak{O}_{(l)}\coprod\mathfrak{O}_{(1)}\times\mbox{Irr}(C_{l})$. The
correspondence is given by assigning to elements
$\\{\rho_{1},\ldots,\rho_{l}\\}$ of $\mathfrak{O}_{(l)}$ the induced
representation $i_{K^{l}}^{C_{l}\wr K}(\rho_{1})$, and on
$\mathfrak{O}_{(l)}\times\mbox{Irr}(C_{l})$ by assigning to elements
$(\rho,e)$ the representation whose underlying $K^{l}$ representation is
$\otimes^{l}\rho$, and whose action by $C$ is given by permuting the tensor
product, and to $(\\{\rho\\},\chi)$ the representation $(\rho,e)\otimes\chi$.
The restriction map $i^{*}$ is now clearly surjective onto the invariant
subring $\mathfrak{R}_{l}[K^{l}]$, which is the first result. The kernel ideal
is clearly spanned by differences of the form
$\langle\rho,\chi\rangle-\langle\rho,\chi_{0}\rangle$, where $\chi_{0}$
denotes the trivial character of $C_{l}$. Such elements lie in
$I_{\mathfrak{R}_{l}[C_{l}]}\cdot\mathfrak{R}_{l}[C_{l}\wr K]$, where
$\mathfrak{R}_{l}[C_{l}\wr K]$ is an $\mathfrak{R}_{l}[C_{l}]$-algebra via the
homomorphism induced by the group homomorphism $C_{l}\wr K\rightarrow C_{l}$.
Since the augmentation ideal in $\mathfrak{R}_{l}[C_{l}]$ has trivial $l$-th
power, the second result follows. The extension to profinite groups is a
straightforward direct limit argument. $\square$
Consider an $l$-Sylow subgroup $\mathbb{S}(l,n)$ of $\Sigma_{l^{n}}$. It is
well-known that this group can be described as an iterated wreath product
$\mathbb{S}(l,n)=\wr^{i}C_{l}$
where $C_{l}$ denotes a cyclic group of order $l$. We also consider the
semidirect product
$\mathfrak{S}(l,n)=\mathbb{S}(l,n)\ltimes\mathbb{Z}_{l}^{l^{n}}$.
###### Proposition 6.2
The image of the homomorphism
$i^{*}:\mathfrak{R}_{l}[\mathfrak{S}(l,n)]\rightarrow\mathfrak{R}_{l}[\mathbb{Z}_{l}^{l^{n}}]$
is the invariant subring
$\mathfrak{R}_{l}[\mathbb{Z}_{l}^{l^{n}}]^{\mathbb{S}(l,n)}$. If we let
$\mathfrak{I}$ denote the kernel of the homomorphism, then
$\mathfrak{I}^{l^{n}}=\\{0\\}$.
Proof: We first note that $\mathfrak{S}(l,n)\cong
C_{l}\wr\mathfrak{S}(l,n-1)$. We now use this observation to perform an
induction over $n$, as follows. The $n=1$ case is Proposition 6.1. For the
case $n\geq 1$, the observation above show us the restriction map
$\mathfrak{R}_{l}[\mathfrak{S}(l,n)]\rightarrow\mathfrak{R}_{l}[\mathfrak{S}(l,n-1)^{l}]\cong\otimes^{l}\mathfrak{R}_{l}[\mathfrak{S}(l,n-1)]$
surjects on to the invariant subring
$(\otimes^{l}\mathfrak{R}_{l}[\mathfrak{S}(l,n-1)])^{C_{l}}$. We observe the
general fact that if $f:V\rightarrow W$ is a homomorphism of
$\mathbb{F}_{l}$-vector spaces, then the map
$\otimes^{l}f:(\otimes^{l}V)^{C_{l}}\rightarrow(\otimes^{l}W)^{C_{l}}$
is also surjective, since the map $f$ admits a right inverse. By the inductive
hypothesis, the map
$\mathfrak{R}_{l}[\mathfrak{S}(l,n-1)]\rightarrow\mathfrak{R}_{l}[\mathbb{Z}^{l^{n-1}}_{l}]$
surjects onto the ring of invariants under the $\mathfrak{S}(l,n-1)$ action on
$\mathfrak{R}_{l}[\mathbb{Z}^{l^{n-1}}_{l}]$. What the induction now shows is
that the image of the restriction map
$\mathfrak{R}_{l}[\mathfrak{S}(l,n)]\rightarrow\mathfrak{R}_{l}[\mathbb{Z}^{l^{n}}_{l}]$
surjects onto the $C_{l}$ invariant part of
$(\otimes^{l}(\mathfrak{R}_{l}[\mathbb{Z}^{l^{n-1}}_{l}])^{C_{l}}$. Now it is
easy to check that this invariant part is exactly identified with the
$\mathfrak{S}(l,n)$-invariant part of
$\mathfrak{R}_{l}[\mathbb{Z}_{l}^{l^{n}}]$ under the identification
$\mathfrak{R}_{l}[\mathbb{Z}_{l}^{l^{n}}]\cong\otimes^{l}\mathfrak{R}_{l}[\mathbb{Z}_{l}^{l^{n-1}}]$
which concludes the proof, which proves the first assertion. The second is an
easy induction using Proposition 6.1. $\square$
We will need to work with some composite order finite groups as well. We will
need some terminology.
###### Definition 6.1
For any finite group $G$, we will write $\hat{\mathfrak{R}}_{l}[G]$ for the
completion of the ring $\mathfrak{R}_{l}[G]$ at its augmentation ideal. For
any profinite group $G$ and finite quotient $Q$ of $G$, we define
$\hat{\mathfrak{R}}_{l}[G;Q]=\hat{\mathfrak{R}}_{l}[Q]\mathrel{\mathop{\otimes}\limits_{\mathfrak{R}_{l}[Q]}}\mathfrak{R}_{l}[G]$
We collect some elementary properties.
###### Proposition 6.3
Let $G$ be a finite group.
1. 1.
If $G$ is an $l$-group, then the natural homomorphism
$\mathfrak{R}_{l}[G]\rightarrow\hat{\mathfrak{R}}_{l}[G]$ is an isomorphism.
More generally, let $\Gamma$ denote a profinite group, with a quotient $G$
which is a finite $l$-group, then the natural homomorphism
$\mathfrak{R}_{l}[G]\rightarrow\hat{\mathfrak{R}}_{l}[\Gamma;G]$ is an
isomorphism.
2. 2.
For any subgroup $H\subseteq G$, the induction homomorphism
$i_{H}^{G}:R[H]\rightarrow R[G]$ naturally extends to a homomorphism
$\hat{i}_{H}^{G}:\hat{\mathfrak{R}}_{l}[H]\rightarrow\hat{\mathfrak{R}}_{l}[G]$.
3. 3.
We let $G_{l}$ denote the $l$-Sylow subgroup of $G$. Then the restriction map
$\hat{\mathfrak{R}}_{l}[G]\rightarrow\hat{\mathfrak{R}}_{l}[G_{l}]$ is an
injection. More generally, let $\Gamma$ denote a profinite group, with a
finite quotient $G$. Let $G_{l}\subseteq G$ be a $l$-Sylow subgroup of $G$.
Then the natural map
$\hat{\mathfrak{R}}_{l}[\Gamma;G]\rightarrow\hat{\mathfrak{R}}_{l}[\Gamma_{l},G_{l}]$
is an injection, where $\Gamma_{l}$ denotes the inverse image of $G_{l}$ in
$\Gamma$ under the projection $\Gamma\rightarrow G$.
4. 4.
For any finite group $G$, the kernel of the (surjective) homomorphism
$R[G]\rightarrow\hat{\mathfrak{R}}_{l}[G]$
is the ideal
$(l)+((l)+I_{G})^{\infty}$
where for any ideal $J^{\infty}$ denotes the intersection $\bigcap_{n}J^{n}$.
More generally, given a profinite group $\Gamma$ and a finite quotient $G$,
the kernel of the homomorphism
$R[\Gamma]\rightarrow\hat{\mathfrak{R}}_{l}[\Gamma;G]$
is the ideal $((l)+((l)+I_{G})^{\infty})R[\Gamma]$.
5. 5.
For any finite $l$-group $G$, the ideal $((l)+I_{G})^{\infty}\subseteq R[G]$
is the zero ideal.
Proof: Part (1) follows from Proposition 3.5 of [34]. For part (2), we first
quote Corollary 3.9 of [34], which asserts that the $I_{H}$ topology and
$I_{G}$ topology on $R[H]$ are the same. It is easy to check that this extend
directly to $\mathfrak{R}_{l}$. To prove that we have a natural extension of
$i_{H}^{G}$, we need to check that for every $n$, there is a $d(n)$ so that
$i_{H}^{G}(I_{H}^{n(d)})\subseteq I_{G}^{n}$. The observation above now shows
that it suffices to prove that for every $n$, there is some number $N(n)$ so
that $i_{H}^{G}(I_{G}^{N(n)}\mathfrak{R}_{l}[H])\subseteq
I_{G}^{n}\mathfrak{R}_{l}[G]$. But Frobenius induction has the property that
$i_{H}^{G}(I_{G}^{n}\mathfrak{R}_{l}[H])\subseteq
I_{G}^{n}\mathfrak{R}_{l}[G]$, so we are done. For part (3), we find that the
composite
$\hat{\mathfrak{R}}_{l}[G]\rightarrow\hat{\mathfrak{R}}_{l}[H]\stackrel{{\scriptstyle\hat{i}_{H}^{G}}}{{\rightarrow}}\hat{\mathfrak{R}}_{l}[G]$
has the property that $1$ is carried to an element $\lambda$ whose
augmentation is equal to $\\#(G/G_{l})$. This number is prime to $l$, and so
is a unit in $\mathbb{F}_{l}$. Since the augmentation ideal in
$\hat{\mathfrak{R}}_{l}[G]$ is nilpotent, it follows that $\lambda$ is a unit,
so $\hat{\mathfrak{R}}_{l}[G]$ is in fact a summand (as
$\hat{\mathfrak{R}}_{l}[G]$-modules) of $\hat{\mathfrak{R}}_{l}[H]$. Part (4)
follows immediately from the definitions. Part (5) Is a direct consequence of
Part (1). $\square$
###### Corollary 6.1
The ring homomorphism
$\hat{\mathfrak{R}}_{l}[\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n};\Sigma_{n}]\rightarrow\mathfrak{R}_{l}[\mathbb{Z}_{l}^{n}]$
has image the invariant subring
$\mathfrak{R}_{l}[\mathbb{Z}_{l}^{n}]^{\Sigma_{n}}$. Further, its kernel
$\mathfrak{I}$ satisfies $\mathfrak{I}^{M}=\\{0\\}$, where $M$ is the smallest
power of $l$ which is greater than $n$.
Proof: Let $\Sigma_{n}(l)\subseteq\Sigma_{n}$ denote an $l$-Sylow subgroup.
Letting $\\{\sigma_{1},\ldots,\sigma_{t}\\}$ denote a complete set of coset
representatives for $\Sigma_{n}(l)$ in $\Sigma_{n}$, we have a homomorphism
$\tau:\mathfrak{R}_{l}[\mathbb{Z}_{l}^{n}]^{\Sigma_{n}(l)}\rightarrow\mathfrak{R}_{l}[\mathbb{Z}_{l}^{n}]^{\Sigma_{n}}$
defined by
$\tau(x)=\sum_{i}\sigma_{i}(x)$
The composite
$\mathfrak{R}_{l}[\mathbb{Z}_{l}^{n}]^{\Sigma_{n}}\stackrel{{\scriptstyle
i}}{{\rightarrow}}\mathfrak{R}_{l}[\mathbb{Z}_{l}^{n}]^{\Sigma_{n}(l)}\stackrel{{\scriptstyle\tau}}{{\rightarrow}}\mathfrak{R}_{l}[\mathbb{Z}_{l}^{n}]^{\Sigma_{n}}$
is equal to multiplication by the index of $\Sigma_{n}(l)$ in $\Sigma_{n}$,
from which it follows that $\tau$ is a surjection onto a summand. It is now
easy to check that there is a commutative diagram
$\begin{diagram}$
where the left hand vertical arrow is induction from
$\Sigma_{n}(l)\ltimes\mathbb{Z}_{l}^{n}$ to
$\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n}$. The upper horizontal arrow is
surjective by Proposition 6.2, and we have observed that $\tau$ is surjective.
It follows that the lower horizontal arrow is also surjective, which is the
first assertion in the Corollary. The second statement follows directly from
(3) in Proposition 6.3. $\square$
###### Corollary 6.2
Let $n$ be a positive integer, and let $s$ and $t$ be positive integers so
that $s+t=n$. The the homomorphism
$\hat{\mathfrak{R}}_{l}[(\Sigma_{s}\times\Sigma_{t})\ltimes\mathbb{Z}_{l}^{n}]\rightarrow\mathfrak{R}_{l}[\mathbb{Z}_{l}^{n}]$
has image equal to the invariant subring
$\mathfrak{R}_{l}[\mathbb{Z}_{l}^{n}]^{\Sigma_{s}\times\Sigma_{t}}$, and there
is an integer $M$ so that the $M$-th power of the kernel ideal is $=\\{0\\}$.
Proof: Immediate from Corollary 6.1, since this example is simply a tensor
product of two of the examples to which it applies. $\square$
###### Corollary 6.3
Let $n,s$, and $t$ be as in the preceding corollary. Then
$\hat{\mathfrak{R}}_{l}[(\Sigma_{s}\times\Sigma_{t})\ltimes\mathbb{Z}_{l}^{n};\Sigma_{s}\times\Sigma_{t}]$
becomes an
$\hat{\mathfrak{R}}_{l}[\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n};\Sigma_{n}]$-algebra
under the ring homomorphism induced by the inclusion
$(\Sigma_{s}\times\Sigma_{t})\ltimes\mathbb{Z}_{l}^{n}\rightarrow\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n}$.
We let $\hat{I}[-]\hookrightarrow\hat{\mathfrak{R}}_{l}[-]$ denote the
augmentation ideal. Then there is an integer $M$ so that
$\hat{I}[(\Sigma_{s}\times\Sigma_{t})\ltimes\mathbb{Z}_{l}^{n};\Sigma_{s}\times\Sigma_{t}]^{M}\subseteq\hat{I}[\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n};\Sigma_{n}]\cdot\hat{\mathfrak{R}}_{l}[(\Sigma_{s}\times\Sigma_{t})\ltimes\mathbb{Z}_{l}^{n};\Sigma_{s}\times\Sigma_{t}]$
Proof: This is a direct application of Theorem 6.2, with
$R=\hat{\mathfrak{R}}_{l}[(\Sigma_{s}\times\Sigma_{t})\ltimes\mathbb{Z}_{l}^{n};\Sigma_{s}\times\Sigma_{t}]$,
$R^{\prime}=\hat{\mathfrak{R}}_{l}[\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n};\Sigma_{n}]$,
$A=\mathfrak{R}_{l}[\mathbb{Z}_{l}]$, $f$ is the homomorphism induced by the
inclusion
$(\Sigma_{s}\times\Sigma_{t})\ltimes\mathbb{Z}_{l}^{n}\hookrightarrow\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n}$,
and where $\pi$ and $\pi^{\prime}$ are the homomorphisms induced by the
inclusions
$\mathbb{Z}_{l}^{n}\hookrightarrow(\Sigma_{s}\times\Sigma_{t})\ltimes\mathbb{Z}_{l}^{n}$
and $\mathbb{Z}_{l}^{n}\hookrightarrow\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n}$
respectively. That the hypothesis of Theorem 6.2 is satisfied is guaranteed by
Corollary 6.2. $\square$
###### Corollary 6.4
Fix $n,s$, and $t$ as in the preceding corollary. Then there is an integer $M$
so that
$I_{(\Sigma_{s}\times\Sigma_{t})\ltimes\mathbb{Z}_{l}^{n}}^{M}\subseteq((l)+I_{\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n}}+((l)+I_{\Sigma_{s}\times\Sigma_{t}})^{\infty})R[\Sigma_{s}\times\Sigma_{t}\ltimes\mathbb{Z}_{l}^{n}]$
Proof: Immediate consequence of Part (4) of Proposition 6.3, together with
Corollary 6.3. $\square$
## 7 Equivariant $K$-theory of ${\cal A}_{\eta}$ and ${\cal T}_{\eta}$
We wish to study the map $K^{G}{\cal A}_{\eta}\rightarrow K^{G}{\cal
T}_{\eta}$, with an eye to proving that it induces an equivalence on
completions along the homomorphism of commutative ring spectra
$\mathfrak{K}^{G}(k)\rightarrow\mathbb{H}_{l}$. More generally, we will want
to study the map $K^{G}({\cal A}_{\eta}\times W)\rightarrow K^{G}({\cal
T}_{\eta}\times W)$, and prove that it induces an equivalence on completions,
where $W$ is a Noetherian approximable Affine scheme with continuous
$G$-action. We will carry out the analysis in the case $W=Spec(k)$, and then
indicate any necessary modifications in the general case. In order to do this,
we will create a filtration by abelian subcategories on the category of all
finitely generated modules over affine space whose support consists of
coordinate hyperplanes (including the full affine space). Of course, this
filtration will induce a filtration of the $K$-theory spectra associated to
the categories, via application of the localization sequence. This section is
entirely devoted to the analysis of the subquotients in this filtration as
modules over the weak ring spectra $K(k)$ and $K^{G}(k)$. We emphasize that
all calculations in this section involve only modules over weak ring spectra.
These calculations will then be used in the next section to draw conclusions
about the derived completions of modules over $\mathfrak{K}^{G}(k)$, via
comparisons such as the ones performed in [8], section 4.
Let $\mathbb{A}=\mathbb{A}(k)=Spec(k[x_{1},x_{2},\ldots,x_{n}])$ be the
$n$-dimensional affine space over $k$, with explicit choice of coordinates
$x_{i}$. Of course, Theorem 7 of [31] tells us that $K(\mathbb{A})\cong K(k)$.
For any $1\leq s\leq n$, we let $V(s)$ denote the hyperplane defined by
$x_{s}=0$. More generally, let $S\subseteq\\{x_{1},\ldots,x_{n}\\}$, and
define
$V(S)=\bigcap_{s\in S}V(s)$
For any $i$, we will let $V_{i}$ denote the union
$V_{i}=\bigcup_{\\#(S)\geq i}V(S)$
Note that $V_{i+1}\subseteq V_{i}$. We will also denote by $\mathbb{T}$ the
scheme $\mathbb{A}-\cup_{s}V_{s}$, which is $Spec(k[x_{1}^{\pm
1},\ldots,x_{n}^{\pm 1}]$. We will refer to subvarieties of the form $V(S)$ as
coordinate subspaces, and the corresponding ideals as coordinate ideals. Note
that the case $S=\emptyset$ is included, with the full $\mathbb{A}$ as
corresponding variety.
###### Definition 7.1
Let $\mbox{Mod}_{i}(\mathbb{A})$ to be the subcategory of finitely generated
modules whose support is contained in $V_{i}$. For a set
$S\subseteq\\{1,\ldots,n\\}$ of cardinality $i$, we also define
$\mbox{Mod}_{i}(\mathbb{A};S)\subseteq\mbox{Mod}_{i}(\mathbb{A})$ to be the
subcategory of modules supported on the subvariety $V(S)\cup V_{i+1}$. We note
that all the categories $\mbox{Mod}_{i}(\mathbb{A})$ and
$\mbox{Mod}_{i}(\mathbb{A};S)$ are abelian categories, and all inclusions are
inclusions of Serre subcategories. Let $\mathbb{A}_{S}$ denote the subscheme
$\mathbb{A}-\bigcup_{s\notin S}V(s)$, and
$\overline{\mathbb{A}}_{S}\subseteq\mathbb{A}_{S}$ denote the subscheme
$\\{(x_{1},\ldots,x_{n})|x_{i}=0\mbox{ for }i\in S\\}$. We also let $\mbox{\em
Nil}(\mathbb{A};S)$ denote the category of coherent $\mathbb{A}_{S}$-modules
on which the generators $\\{x_{i}\\}_{i\in S}$ act nilpotently.
We clearly have an increasing sequence of subcategories
$\\{0\\}=\mbox{Mod}_{n+1}(\mathbb{A})\subseteq\mbox{Mod}_{n}(\mathbb{A})\subseteq\mbox{Mod}_{n-1}(\mathbb{A})\subseteq\cdots\subseteq\mbox{Mod}_{1}(\mathbb{A})\subseteq\mbox{Mod}_{0}(\mathbb{A})=\mbox{Mod}(\mathbb{A})$
###### Proposition 7.1
There is a natural functor from the quotient abelian category $\mbox{\em
Mod}(\mathbb{A})/\mbox{\em Mod}_{1}(\mathbb{A})$ to the category $\mbox{\em
Mod}(\mathbb{T})$, which induces an equivalence of $K$-theory spectra.
Proof: This is simply the localization sequence associated to the removal of
$V_{1}$ from $V_{0}\ \mathbb{A}$, which is $\mathbb{T}$. $\square$
We let $\Phi(\mathbb{A},\mathbb{T})$ denote the homotopy fiber of the map of
$K$-theory spectra $K(\mathbb{A})\rightarrow K(\mathbb{T})$. It follows
readily from Proposition 7.1 above that $\Phi(\mathbb{A},\mathbb{T})$ is
equivalent to the $K$-theory spectrum $K(\mbox{Mod}_{1})$.
###### Proposition 7.2
There is a filtration of $\Phi=\Phi(\mathbb{A},\mathbb{T})$ by spectra
$\Phi_{i}=\Phi_{i}(\mathbb{A},\mathbb{T})$, with
$*\simeq\Phi_{n+1}\subseteq\Phi_{n}\subseteq\Phi_{n-1}\subseteq\cdots\subseteq\Phi_{2}\subseteq\Phi_{1}=\Phi$
and so that $\Phi_{i}/\Phi_{i+1}\cong K(\mbox{\em
Mod}_{i}(\mathbb{A})/\mbox{\em Mod}_{i+1}(\mathbb{A}))$.
Proof: The spectra $\Phi_{i}$ can be taken to be the homotopy fibers of the
maps of spectra
$K(\mathbb{A})\rightarrow
K(\mbox{Mod}(\mathbb{A})/\mbox{Mod}_{i}(\mathbb{A}))$
$\square$
Each inclusion
$\mbox{Mod}_{i+1}(\mathbb{A})\hookrightarrow\mbox{Mod}_{i}(\mathbb{A})$ is an
inclusion of an abelian subcategory, and it is clear from the discussion of
quotient abelian categories in [38] that the quotient
$\mbox{Mod}_{i}(\mathbb{A})/\mbox{Mod}_{i+1}(\mathbb{A})$ can be analyzed as
follows. For each $S$ with $\\#(S)=i$, we may consider the inclusion
$i_{S}:\mbox{Mod}_{i}(\mathbb{A};S)\hookrightarrow\mbox{Mod}_{i}(\mathbb{A})$,
and therefore the inclusion of subquotients
$j_{S}:\mbox{Mod}_{i}(\mathbb{A};S)/\mbox{Mod}_{i+1}(\mathbb{A})\hookrightarrow\mbox{Mod}_{i}(\mathbb{A})/\mbox{Mod}_{i+1}(\mathbb{A})$.
We may then sum over all $S$ of cardinality $i$ and apply $K$-theory to obtain
$\bigvee_{\\#(S)=i}j_{S}:\bigvee_{\\#(S)=i}K(\mbox{Mod}_{i}(\mathbb{A};S)/\mbox{Mod}_{i+1}(\mathbb{A}))\longrightarrow
K(\mbox{Mod}_{i}(\mathbb{A})/\mbox{Mod}_{i+1}(\mathbb{A}))$ (7–6)
This functor is an equivalence since as in the proof of Theorem 5.4 of [31] it
can be checked that the underlying functor is an equivalence of categories.
Further, there is a natural equivalence of abelian categories
$\mbox{Mod}_{i}(\mathbb{A};S)/\mbox{Mod}_{i+1}(\mathbb{A})\simeq\mbox{Nil}(\mathbb{A};S)$
and hence equivalences of $K$-theory spectra of these categories
$K(\mbox{Mod}_{i}(\mathbb{A};S)/\mbox{Mod}_{i+1}(\mathbb{A}))\simeq
K(\mbox{Nil}(\mathbb{A};S))\simeq K(\mbox{Mod}(\overline{\mathbb{A}}_{{S}}))$
The second of the equivalences is proved as in [31] by the devissage
principle.
We will also consider the situation of a product of a fixed Noetherian affine
scheme $W$ with $\mathbb{A}$ and $\mathbb{T}$. In this case, we define the
subcategory $\mbox{Mod}_{i}(\mathbb{A},W)\subseteq\mbox{Mod}(\mathbb{A}\times
W)$ to consist of all $\mathbb{A}\times W$-modules whose support is contained
in $\pi^{-1}(V_{i})$, where $\pi:\mathbb{A}\times W\rightarrow\mathbb{A}$ is
the projection. There are obvious analogues $\mbox{Mod}_{i}(\mathbb{A},W;S)$
and $\mbox{Nil}(\mathbb{A},W;S)$ of $\mbox{Mod}_{i}(A;S)$ and
$\mbox{Nil}(\mathbb{A};S)$.
It is further clear from Section 2.2 that all the spectra
$K(\mbox{Mod}_{i}(\mathbb{A},W)),K(\mbox{Mod}_{i}(\mathbb{A},W;S)),\mbox{ and
}K(\mbox{Nil}(\mathbb{A},W;S))$
are module spectra over the weak ring spectrum $K(\mathbb{A}\times W)$, that
all the inclusion and restriction functors induce maps of $K(\mathbb{A}\times
W)$-module spectra, and therefore that the quotient spectra
$K(\mbox{Mod}_{i}(\mathbb{A},W))/K(\mbox{Mod}_{i+1}(\mathbb{A},W))$ are also
$K(\mathbb{A}\times W)$-module spectra. It is also clear that
$K(\mbox{Nil}(\mathbb{A},W;S))$ is naturally a $K(\mathbb{A}_{S}\times
W)$-module spectrum, in a way which extends the $K(\mathbb{A}\times W)$-module
structure via the map of weak ring spectra $K(\mathbb{A}\times W)\rightarrow
K(\mathbb{A}_{S}\times W)$. Moreover, it is clear that the decomposition of
spectra
$K(\mbox{Mod}_{i}(\mathbb{A},W))/K(\mbox{Mod}_{i+1}(\mathbb{A},W))\cong\bigvee_{\\#(S)=i}K(\mbox{Nil}(\mathbb{A},W;S))\cong\bigvee_{\\#(S)=i}K(\mbox{Mod}(\overline{\mathbb{A}}_{{S}}\times
W))$
given by the homotopy inverse to the map $\bigvee_{\\#(S)=i}K(j_{s})$ from 7–6
above is also an equivalence of $K(\mathbb{A}\times W)$-module spectra, where
the $K(\mathbb{A}\times W)$ action on the summand
$K(\mbox{Nil}(\mathbb{A},W;S))$ and
$K(\mbox{Mod}(\overline{\mathbb{A}}_{{S}}\times W))$ are obtained by
restriction of scalars along the map of weak ring spectra $K(\mathbb{A}\times
W)\rightarrow K(\mathbb{A}_{S}\times W)$.
We next begin the analysis of how this filtration works equivariantly. Let
$\rho$ be any monomial $k$-linear representation of a finite group $G$, and
let $\\{x_{1},\ldots,x_{n}\\}$ denote any basis for the representation space
$V_{\rho}$ in which the representation is monomial. Let
$\overline{\rho}:G\rightarrow\Sigma_{n}$ denote the corresponding permutation
representation. As usual, we write $\mathbb{A}=\mathbb{A}(k,\rho)$ for the
associated affine variety, and therefore on the associated module category. It
is clear that the filtration of $\mbox{Mod}(\mathbb{A})$ by the subcategories
$\mbox{Mod}_{i}(\mathbb{A})$ is invariant under the action defined by $\rho$,
and we therefore easily obtain a filtration of $\mbox{Mod}^{G}(\mathbb{A})$ by
Serre subcategories $\mbox{Mod}_{i}^{G}(\mathbb{A})$. The equivariant version
of Proposition 7.2 holds. Let $\Phi^{G}=\Phi^{G}(\mathbb{A},\mathbb{T})$
denote the homotopy fiber of the map of spectra
$K^{G}(\mathbb{A})\rightarrow K^{G}(\mathbb{T})$
###### Proposition 7.3
There is a filtration of $\Phi^{G}$ by spectra
$\Phi_{i}^{G}=\Phi_{i}^{G}(\mathbb{A},\mathbb{T})$, with
$*\simeq\Phi_{n+1}^{G}\subseteq\Phi_{n}^{G}\subseteq\Phi_{n-1}^{G}\subseteq\cdots\subseteq\Phi_{2}^{G}\subseteq\Phi_{1}^{G}=\Phi^{G}$
and so that $\Phi^{G}_{i}/\Phi^{G}_{i+1}\cong K^{G}(\mbox{\em
Mod}_{i}(\mathbb{A})/\mbox{\em Mod}_{i+1}(\mathbb{A}))$.
We will have several categories of $\mathbb{A}$-modules which are closed under
the action of $G$ or of subgroups of $G$. For each such category, we can then
denote its equivariant version with the superscript $G$, and define its
equivariant $K$-theory as the $K$-theory of the equivariant version of the
category. For instance, we will write $K^{G}(\mbox{Mod}_{i}(\mathbb{A}))$ for
$K(\mbox{Mod}^{G}_{i}(\mathbb{A}))$. We now identify the subquotients
$\mbox{Mod}_{i}^{G}(\mathbb{A})/\mbox{Mod}_{i+1}^{G}(\mathbb{A})$. Let ${\cal
P}_{i}$ denote the set of subsets of $\\{1,\ldots,n\\}$ of cardinality $i$,
and let
${\cal P}_{i}=\mathfrak{o}_{1}\cup\cdots\cup\mathfrak{o}_{t}$
be an orbit decomposition, with orbit representatives
$S_{j}\in\mathfrak{o}_{j}$. Note that for each $j$, the subcategory
$\mbox{Mod}_{i}(\mathbb{A};\mathfrak{o}_{j})$ is invariant under the action of
$G$. Then an argument entirely analogous to the non-equivariant one above
shows that
$\mbox{Mod}_{i}^{G}(\mathbb{A})/\mbox{Mod}_{i+1}^{G}(\mathbb{A})\cong\bigvee_{j=1}^{t}\mbox{Mod}^{G}_{i}(\mathbb{A};\mathfrak{o}_{j})/\mbox{Mod}^{G}_{i+1}(\mathbb{A})$
Moreover, it is now easy to see that each summand is described by
$\mbox{Mod}^{G}_{i}(\mathbb{A};\mathfrak{o}_{j})/\mbox{Mod}^{G}_{i+1}(\mathbb{A})\cong\mbox{Mod}^{G_{j}}_{i}(\mathbb{A};S_{j})/\mbox{Mod}^{G_{j}}_{i+1}(\mathbb{A})$
where $G_{j}$ is the stabilizer of the orbit representative $S_{j}$, and that
further we have an equivalence of abelian categories
$\mbox{Mod}^{G_{j}}_{i}(\mathbb{A},S_{j})/\mbox{Mod}_{i+1}^{G_{j}}(\mathbb{A})\cong\mbox{Nil}^{G_{j}}(\mathbb{A};{{S}_{j}})$
(7–7)
Also, devissage gives us an equivalence of spectra
$K(\mbox{Nil}^{G_{j}}(\mathbb{A};S_{j}))\cong
K^{G_{j}}(\mbox{Mod}(\overline{\mathbb{A}}_{{S}_{j}}))$ (7–8)
and therefore a decomposition of $K$-theory spectra
$K(\mbox{Mod}^{G}_{i}(\mathbb{A}))/K(\mbox{Mod}^{G}_{i+1}(\mathbb{A}))\cong\bigvee_{j=1}^{t}K^{G_{j}}(\mbox{Nil}(\mathbb{A};{{S}_{l}}))\cong\bigvee_{j=1}^{t}K^{G_{j}}(\overline{\mathbb{A}}_{{S}_{j}})$
(7–9)
The corresponding decomposition for the case of a product with a fixed affine
Noetherian scheme with $G$-action $W$
$K(\mbox{Mod}^{G}_{i}(\mathbb{A},W))/K(\mbox{Mod}^{G}_{i+1}(\mathbb{A},W))\cong\bigvee_{j=1}^{t}K^{G_{j}}(\mbox{Nil}(\mathbb{A},W;S_{l}))\cong\bigvee_{j=1}^{t}K^{G_{j}}(\overline{\mathbb{A}}_{{S}_{j}}\times
W)$ (7–10)
follows similarly
We now let $G$ be a profinite group, and let $\eta$ denote a fixed $l$-BAR of
$G$. We also let $W$ denote a continuous Noetherian affine $G$-scheme. In
practice, we will be concerned only with the case where $W=Spec(k)$, and the
$G$-action is trivial, or the case where $G=G_{F}$ is the absolute Galois
group of an extension $F$ of $k$, and $W=Spec(\overline{F})$. We wish to
analyze $K^{G}({\cal A}_{\eta}\times W)$, $K^{G}({\cal T}_{\eta}\times W)$,
and the map $K^{G}({\cal A}_{\eta}\times W)\rightarrow K^{G}({\cal
T}_{\eta}\times W)$. We obtain $k$-linear representations $\rho_{s}\cdot\eta$
of $G$, and the corresponding varieties ${\cal A}_{\eta}^{s}$ and ${\cal
T}_{\eta}^{s}$. We also denote the affine coordinate ring of ${\cal
A}_{\eta}^{s}$ by $A_{s}$. We also defined the $l$-th power map $\theta:{\cal
A}_{\eta}^{s+1}\rightarrow{\cal A}_{\eta}^{s}$, which is equivariant.Because
the map $\theta$ respects the coordinate axes, it is clear that it carries
coordinate ideals into coordinate ideals. From the definition of the
subcategories $\mbox{Mod}^{G}_{i}$, it follows readily that the functor
$A_{s+1}\mathrel{\mathop{\otimes}\limits_{A_{s}}}-$ carries
$\mbox{Mod}_{i}^{G}({\cal A}_{\eta}^{s})\subseteq\mbox{Mod}^{G}({\cal
A}_{\eta}^{s})$ to the subcategory $\mbox{Mod}_{i}^{G}({\cal
A}_{\eta}^{s+1})\subseteq\mbox{Mod}^{G}({\cal A}_{\eta}^{s+1})$. Note that
these tensoring functors are exact because ${A}_{s+1}$ is always a flat
$A_{s}$-algebra. This permits us to define categories
$\mbox{Mod}_{i}^{G}({\cal A}_{\eta})$ as the colimits of the systems
$\cdots\rightarrow\mbox{Mod}_{i}^{G}({\cal
A}_{\eta}^{s-1})\rightarrow\mbox{Mod}_{i}^{G}({\cal{A}}_{\eta}^{s})\rightarrow\mbox{Mod}_{i}^{G}({\cal
A}_{\eta}^{s+1})\rightarrow\cdots$
A more intrinsic way to describe the subcategories of $\mbox{Mod}_{i}({\cal
A_{\eta}})$ described above is as follows. For any $s\in\\{1,\ldots,n\\}$, we
let $V(s)$ denote the subscheme defined by the ideal
$(x_{s},x_{s}^{\frac{1}{l}},x_{s}^{\frac{1}{l^{2}}},\ldots)$. We can now
construct $V(S)$ for $S\subseteq\\{1,\ldots,n\\}$ as before, and the
corresponding module categories of finitely presented modules supported on the
corresponding subschemes as well.
We note that the analyses above apply equally well to a situation where we
take products with a fixed affine Noetherian $G$-scheme. The arguments from
above are adapted in a straightforward way.
We now have the analogue of the Proposition 7.3. Let $\Phi^{G}({\cal
A}_{\eta},{\cal T}_{\eta})$ denote the homotopy fiber of the map
$K^{G}({\cal A}_{\eta})\rightarrow K^{G}({\cal T}_{\eta})$
###### Proposition 7.4
There is a filtration of $\Phi^{G}({\cal A}_{\eta},{\cal T}_{\eta})$ by
spectra $\Phi_{i}^{G}({\cal A}_{\eta},{\cal T}_{\eta})$, with
$*\simeq\Phi_{n+1}^{G}({\cal A}_{\eta},{\cal
T}_{\eta})\subseteq\Phi_{n}^{G}({\cal A}_{\eta},{\cal
T}_{\eta})\subseteq\Phi_{n-1}^{G}({\cal A}_{\eta},{\cal
T}_{\eta})\subseteq\cdots\subseteq\Phi_{2}^{G}({\cal A}_{\eta},{\cal
T}_{\eta})\subseteq\Phi_{1}^{G}({\cal A}_{\eta},{\cal
T}_{\eta})=\Phi^{G}({\cal A}_{\eta},{\cal T}_{\eta})$
and so that
$\Phi^{G}_{i}({\cal A}_{\eta},{\cal T}_{\eta})/\Phi^{G}_{i+1}({\cal
A}_{\eta},{\cal T}_{\eta})\cong K^{G}(\mbox{\em Mod}_{i}({\cal
A}_{\eta})/\mbox{\em Mod}_{i+1}({\cal A}_{\eta}))$
Further, for any orbit $\mathfrak{o}_{j}$ under the action of $G$ on ${\cal
P}_{i}$, $A_{s+1}\mathrel{\mathop{\otimes}\limits_{A_{s}}}-$ carries the
subcategory $\mbox{Mod}^{G}_{i}({\cal A}^{s}_{\eta},\mathfrak{o}_{j})$ into
the subcategory $\mbox{Mod}^{G}_{i}({\cal A}_{\eta}^{s+1},\mathfrak{o}_{j})$.
It follows that the decompositions 7–10 are respected by the map $\theta$. We
can therefore define subcategories
$\mbox{Mod}_{i}^{G}({\cal
A}_{\eta};\mathfrak{o})\subseteq\mbox{Mod}_{i}^{G}({\cal A}_{\eta})$
as the colimits of the system of categories
$\cdots\rightarrow\mbox{Mod}_{i}^{G}({\cal
A}_{\eta}^{s-1};\mathfrak{o})\rightarrow\mbox{Mod}_{i}^{G}({\cal
A}^{s}_{\eta};\mathfrak{o})\rightarrow\mbox{Mod}_{i}^{G}({\cal
A}_{\eta}^{s+1};\mathfrak{o})\rightarrow\cdots$
and we obtain a decomposition
$K^{G}(\mbox{Mod}_{i}({\cal A}_{\eta})/\mbox{Mod}_{i+1}({\cal
A}_{\eta}))\cong\bigvee_{j=1,\ldots t}K^{G}(\mbox{Mod}_{i}({\cal
A}_{\eta};\mathfrak{o_{j}})/\mbox{Mod}_{i+1}({\cal A}_{\eta}))$ (7–11)
where $\\{\mathfrak{o}_{i},\ldots,\mathfrak{o}_{t}\\}$ is an orbit
decomposition of the action of $G$ on the subsets of cardinality $i$ of
$\\{1,\ldots,n\\}$.
We will now focus on analyzing the effect of the map
$K(\mbox{Mod}^{G}_{i}({\cal
A}_{\eta}^{s};\mathfrak{o}_{j})/\mbox{Mod}^{G}_{i+1}({\cal
A}_{\eta}^{s}))\rightarrow K(\mbox{Mod}^{G}_{i}({\cal
A}_{\eta}^{s+1};\mathfrak{o}_{j})/\mbox{Mod}^{G}_{i+1}({\cal
A}_{\eta}^{s+1}))$
induced by $\theta$ on the individual factors, so as to understand the groups
$\pi_{*}K^{G}(\mbox{Mod}_{i}({\cal
A}_{\eta};\mathfrak{o_{j}})/\mbox{Mod}_{i+1}({\cal A}_{\eta}))$
We let $S_{j}$ denote a representative for the orbit $\mathfrak{o}_{j}$, and
let $G_{j}\subseteq G$ denote the stabilizer of $S_{j}$. We also have the
identifications in 7–7 and 7–8 above, and we wish to identify the effect of
$\theta$ under this identification. The category $\mbox{Nil}^{G_{j}}({\cal
A}_{\eta}^{s};{{S}_{j}})$ is described as follows. We suppose that an orbit
representative for the set $S_{j}$ is the set $\\{1,\ldots,i\\}$. If not,
perform a reordering of the indexing. The permutation representation of
$G_{j}$ preserves the decomposition
$\\{1,\ldots,i\\}\coprod\\{i+1,\ldots,n\\}$ (7–12)
and yields decompositions $\rho_{s}\cdot\eta\simeq\sigma_{s}\oplus\tau_{s}$
and $\rho_{s+1}\cdot\eta\simeq\sigma_{s+1}\oplus\tau_{s+1}$, and corresponding
decompositions of schemes with $G$-action
${\cal A}_{\eta}^{s}={\cal B}_{\eta}^{s}\times{\cal
C}_{\eta}^{s}=\mathbb{A}(k,\sigma_{s})\times\mathbb{A}(k,\tau_{s})$ (7–13)
for all $s$. We also remove the coordinate hyperplanes corresponding to the
elements of the complement of $S_{j}$, and obtain the variety
$\overline{{\cal A}}^{s}_{\eta}={\cal{B}}^{s}_{\eta}\times\overline{{\cal
C}}^{s}_{\eta}=\mathbb{A}(k,\sigma_{s})\times\mathbb{T}(k,\tau_{s})$
Moreover, the power map $\theta=\theta_{\rho}$ preserves the decomposition, so
that we have a commutative diagram
$\begin{diagram}$
of $G_{j}$-equivariant schemes. The category $\mbox{Nil}^{G_{j}}({\cal
A}_{\eta}^{s};{{S}_{j}})$ corresponds to the category of all finitely
generated $\overline{{\cal A}}_{\eta}^{s}$-modules whose support is
concentrated on the subscheme $\\{0\\}\times\overline{{\cal C}}_{\eta}^{s}$,
and similarly for $\overline{{\cal A}}_{\eta}^{s+1}$. We let
$A_{s}=k[x_{1},\ldots,x_{i},x_{i+1}^{\pm 1},\ldots,x_{n}^{\pm
1}],B_{s}=k[x_{1},\ldots,x_{i}]$, and $C_{s}=k[x_{i+1}^{\pm
1},\ldots,x_{n}^{\pm 1}]$ denote the corresponding affine coordinate rings,
and similarly for $s+1$.
###### Proposition 7.5
Since the product decompositions 7–13 are equivariant, we have the following.
1. 1.
The projections $\overline{\cal A}_{\eta}^{s}\rightarrow\overline{\cal
C}_{\eta}^{s}$ and $\overline{\cal A}_{\eta}^{s+1}\rightarrow\overline{\cal
C}_{\eta}^{s+1}$ induce equivalences $K^{G_{j}}(\overline{\cal
A}_{\eta}^{s})\cong K^{G_{j}}(\overline{\cal C}_{\eta}^{s})$ and
$K^{G_{j}}(\overline{\cal A}_{\eta}^{s+1})\cong K^{G_{j}}(\overline{\cal
C}_{\eta}^{s+1})$ of weak ring spectra.
2. 2.
There are weak equivalences
$\varepsilon_{s}:K^{G_{j}}(\overline{\cal
C}_{\eta}^{s})\xrightarrow{\sim}K(\mbox{\em Nil}^{G_{j}}({\cal
A}_{\eta}^{s};S_{j}))$
and
$\varepsilon_{s+1}:K^{G_{j}}(\overline{\cal
C}_{\eta}^{s+1})\xrightarrow{\sim}K(\mbox{\em Nil}^{G_{j}}({\cal
A}_{\eta}^{s+1};S_{j}))$
of modules over the weak ring spectra $K^{G_{j}}({\cal C}_{\eta}^{s})$ and
$K^{G_{j}}({\cal C}_{\eta}^{s})$, which is given on $\pi_{0}$ by
$K_{0}(\varepsilon_{s})(1)=[k\mathrel{\mathop{\otimes}\limits_{k}}C_{s}]$,
where $k$ is the quotient of $B_{s}$ by the maximal ideal of the origin, where
the group action on $k$ is by the identity, and on $C_{s}$ is given by the
already prescribed action. The description of $\varepsilon_{s+1}$ is
identical.
Proof: The first result follows from the homotopy property in [31]. The second
two are straightforward devissage arguments. $\square$
Next, we note that by the homotopy property of $K$-theory, we have
equivalences $K^{G_{j}}(Spec(k))\rightarrow K^{G_{j}}({\cal B}_{\eta}^{s})$
and $K^{G_{j}}(Spec(k))\rightarrow K^{G_{j}}({\cal B}_{\eta}^{s+1})$ induced
by the projections ${\cal B}_{\eta}^{s}\rightarrow Spec(k)$ and ${\cal
B}_{\eta}^{s+1}\rightarrow Spec(k)$, respectively. These equivalences, taken
together with the map $K^{G_{j}}({\cal B}_{\eta}^{s})\rightarrow
K^{G_{j}}({\cal B}_{\eta}^{s+1})$ induced by $\theta_{\sigma}$, give a map
$\overline{\theta}_{\sigma}:K^{G_{j}}(Spec(k))\rightarrow K^{G_{j}}(Spec(k))$,
which we wish to identify. In order to do this, we consider a monomial
$k$-linear representation $\nu$ of a finite group $G$, with the representation
space $V_{\nu}$ of $G$ equipped with a basis $B_{\nu}$ in which the
representation is monomial. Given any vector space $V$ equipped with a basis
$B$, we may construct the symmetric algebra (over $k$) $S[V]$, and factor out
the ideal generated by the set $\\{b^{n}\\}_{b\in B}$. We denote this vector
space by $T_{n}(V,B)$. Note that the construction only depends on the set of
lines spanned by the individual basis elements, and can hence be defined as a
functor on the category of vector spaces equipped with a family of independent
and spanning lines. This means that in our context, we can construct from the
representation $\nu$ a new representation $T_{n}(\nu)$, which has
$T_{n}(V_{\nu},B_{\nu})$ as its representation space. We can now describe the
map $\overline{\theta}_{\sigma}$ as follows.
###### Proposition 7.6
The map $\overline{\theta}_{\sigma}:K^{G_{j}}(k)\rightarrow K^{G_{j}}(k)$ is
characterized up to homotopy by the following two requirements.
1. 1.
$\overline{\theta}_{\sigma}$ is a map of $K^{G_{j}}(k)$-modules.
2. 2.
$\pi_{0}(\overline{\theta}_{\sigma})(1)=T_{l}(\sigma_{s+1})$
Proof: The first statement is immediate since the map $K^{G_{j}}({\cal
B}_{\eta}^{s})\rightarrow K^{G_{j}}({\cal B}_{\eta}^{s+1})$ is a map of
$K^{G_{j}}(k)$-modules. For the second, we pass to affine coordinate rings. We
write
${\cal B}_{\eta}^{s+1}=Spec(B_{s+1})=Spec(k[x_{1},\ldots,x_{i}])$
with the $G_{j}$-action given by an action on the vector space spanned by
$\\{x_{1},\ldots,x_{i}\\}$, which is monomial with respect to this basis. We
have
${\cal B}_{\eta}^{s}=Spec(B_{s})=Spec(k[x_{1}^{l},\ldots,x_{i}^{l}])$
and the map $\theta_{\sigma}:{\cal B}_{\eta}^{s+1}\rightarrow{\cal
B}_{\eta}^{s}$ is given by the evident inclusion $B_{s}\hookrightarrow
B_{s+1}$. The restriction of the $G_{j}$-action to $B_{s}$ factors through the
homomorphism $G_{j}\rightarrow G_{j}$. We are interested in the behavior of
the functor $B_{s+1}\mathrel{\mathop{\otimes}\limits_{B_{s}}}-$ restricted to
the category $\mbox{Nil}^{G_{j}^{s}}({\cal B}_{\eta}^{s})$, the category of
modules whose support is the origin in the affine space ${\cal B}_{\eta}^{s}$.
A devissage argument shows that a generator for
$K_{0}(\mbox{Nil}^{G_{j}}({\cal B}_{\eta}^{s}))$ is the field $k$, regarded as
a module over $B_{s}$ via the trivial (zero) action of each of the generators
$x_{t}^{l}$ and as a $G_{j}$-representation by declaring the action to be the
identity action. We now consider
$B_{s+1}\mathrel{\mathop{\otimes}\limits_{B_{s}}}k$. It can clearly be
identified with $B_{s+1}/(x_{1}^{l},\ldots,x_{i}^{l})$, which as a
$G_{j}$-representation is isomorphic to $T_{l}(\sigma_{s+1})$. It now follows
from the homotopy property that this element in $K_{0}^{G_{j}}({\cal
B}_{\eta}^{s+1})$ is equal to the element $[T_{l}(\sigma_{s+1})]$. $\square$
###### Corollary 7.1
There is a commutative diagram
$\begin{diagram}$
for all $i$, where $\varepsilon_{s}$ and $\varepsilon_{s+1}$ are isomorphisms.
Proof: This is a simple product calculation, noting that ${\cal
A}^{s}\cong{\cal B}_{\eta}^{s}\times{\cal C}_{\eta}^{s}$. $\square$
###### Corollary 7.2
We let $i$ be an integer between $1$ and $n$, that $\mathfrak{o}_{j}$ is an
orbit of the action of $G$ on the collection of subsets of cardinality $i$ in
$\\{1,\ldots,n\\}$ with orbit representative $S_{j}$, and that $G_{j}$ is the
stabilizer of $S_{j}$. Then we have the following.
1. 1.
The graded $R[G]$-module $K_{*}^{G}(\mbox{ \em Mod}_{i}({\cal
A}_{\eta};\mathfrak{o}_{j})/\mbox{\em Mod }_{i}({\cal A}_{\eta}))$ is
isomorphic to the colimit of the system
$K_{*}^{G_{j}}(\overline{{\cal C}}^{0}_{\eta})\stackrel{{\scriptstyle
T_{l}(\sigma_{1})\cdot
K^{G_{j}}_{*}(\theta_{\tau})}}{{\longrightarrow}}K_{*}^{G_{j}}(\overline{{\cal
C}}^{1}_{\eta})\stackrel{{\scriptstyle T_{l}(\sigma_{2})\cdot
K^{G_{j}}_{*}(\theta_{\tau})}}{{\longrightarrow}}K_{*}^{G_{j}}(\overline{{\cal
C}}^{2}_{\eta})\stackrel{{\scriptstyle T_{l}(\sigma_{3})\cdot
K^{G_{j}}_{*}(\theta_{\tau})}}{{\longrightarrow}}\cdots$ (7–14)
where $\sigma_{s}$ is the $i$-dimensional representation of $G$ on the basis
elements corresponding to the subset $S_{j}$ on the affine space
$\mathbb{A}_{s}$. It follows that it is isomorphic to the colimit of the
system
$K_{*}^{G_{j}}(\overline{{\cal C}}_{\eta})\stackrel{{\scriptstyle
T_{l}(\sigma_{1})}}{{\longrightarrow}}K_{*}^{G_{j}}(\overline{{\cal
C}}_{\eta})\stackrel{{\scriptstyle
T_{l}(\sigma_{2})}}{{\longrightarrow}}K_{*}^{G_{j}}(\overline{{\cal
C}}_{\eta})\stackrel{{\scriptstyle
T_{l}(\sigma_{3})}}{{\longrightarrow}}K_{*}^{G_{j}}(\overline{{\cal
C}}_{\eta})\stackrel{{\scriptstyle
T_{l}(\sigma_{4})}}{{\longrightarrow}}\cdots$ (7–15)
2. 2.
Let $\overline{G}_{j}$ denote the image of $G_{j}$ under $\eta$. Then all of
the elements $T_{l}(\sigma_{s})\in R[G_{j}]$ are in the image of the ring
homomorphism $R[\overline{G}_{j}]\rightarrow R[G_{j}]$.
Proof: For part (1), we observe that the colimit (7–14) above can be written
as a two parameter colimit, of the following form.
$\begin{diagram}$
This colimit may be computed in iterative fashion, first in the vertical
direction and then in the horizontal. The diagram (7–15) above is obtained
from this two parameter diagram by performing the vertical colimits. For part
(2) above, it is clear that the representations $\sigma_{s}$ all factor
through $\overline{G}_{j}$ by definition. $\square$
## 8 Completed equivariant $K$-theory of ${\cal T}_{\eta}$
We suppose that we are given an $l$-BAR $\eta$, and construct the spectra
${\cal A}_{\eta}$ and ${\cal T}_{\eta}$. Also let $W$ be a Noetherian
approximable continuous affine $G$-scheme, as in the previous section. We wish
to prove that the map
$\theta:\mathfrak{K}^{G}(k)\cong\mathfrak{K}^{G}({\cal A}_{\eta}\times
W)\rightarrow\mathfrak{K}^{G}({\cal T}_{\eta}\times W)$
induces an equivalence on derived completions at the homomorphism
$\epsilon:\mathfrak{K}^{G}(k)\rightarrow\mathbb{H}_{l}$
where $\mathbb{H}_{l}$ denotes the mod $l$ Eilenberg spectrum. In order to do
this, it will suffice to prove that the map
$\mathbb{H}_{l}\cong\mathfrak{K}^{G}({\cal A}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l}\xrightarrow{\theta\wedge
id}\mathfrak{K}^{G}({\cal T}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l}$
is an equivalence of spectra by Proposition 2.20. This analysis will require
the filtrations discussed in Section 7. The computations there were
exclusively using the weak ring spectra and modules constructed from
Waldhausen’s $S_{.}$-construction. In order to apply those results, we have to
construct the $\mathfrak{K}^{G}$-counterparts to the spectra
$K^{G}(\mbox{Mod}_{i}(-))$ and
$K^{G}(\mbox{Mod}^{G}_{i}(-)/\mbox{Mod}^{G}_{i+1}(-))$. To begin, we suppose
that we are given a monomial representation $l$-BAR of $G$ , and the
corresponding actions of $G$ on ${\cal A}_{\eta}^{s}$ and ${\cal
T}_{\eta}^{s}$. In Section 7, we defined the abelian subcategories
$\mbox{Mod}^{G}_{i}({\cal A}_{\eta}\times W)$ of the category
$\mbox{Mod}^{G}({\cal A}_{\eta}\times W)$. We now define abelian categories
$Q_{i}$ to be the quotient abelian categories $\mbox{Mod}^{G}({\cal
A}_{\eta}\times W)/\mbox{Mod}^{G}_{i}({\cal A}_{\eta}\times W)$. Each of the
categories $Q_{i}$ admits a strictly associative and coherently commutative
tensor produce, and so is therefore suitable input to the Elmendorf-Mandell
machine. Consequently, we may construct commutative
$\mathfrak{K}^{G}(k)$-algebras $\mathfrak{K}^{G}(Q_{i})$, and we have natural
homomorphisms $\mathfrak{K}^{G}(Q_{i})\rightarrow\mathfrak{K}^{G}(Q_{i-1})$ of
such algebras. We now define the module $\mathfrak{M}^{G}_{i}({\cal
A}_{\eta}\times W)$ over $\mathfrak{K}^{G}(k)$ to be the homotopy fiber of the
map
$\mathfrak{K}^{G}(k)({\cal A}_{\eta}\times
W)\rightarrow\mathfrak{K}^{G}(Q_{i})$
The localization theorem in Waldhausen theory shows that
$\mathfrak{M}^{G}_{i}({\cal A}_{\eta}\times W)$ has the same homotopy type as
$K(\mbox{Mod}_{i}^{G}({\cal A}_{\eta}\times W))$ as a weak module spectrum
over $K^{G}({\cal A}_{\eta}\times W)$. We now have natural homomorphisms of
$\mathfrak{K}^{G}(k)$-modules $\mathfrak{M}^{G}_{i+1}({\cal A}_{\eta}\times
W)\rightarrow\mathfrak{M}^{G}_{i}({\cal A}_{\eta}\times W)$, and we can
construct their cofibers $\mathfrak{M}^{G}_{i}({\cal A}_{\eta}\times
W)/\mathfrak{M}^{G}_{i+1}({\cal A}_{\eta}\times W)$ to obtain
$\mathfrak{K}^{G}$ counterparts to the subquotients
$K^{G}(\mbox{Mod}_{i}^{G}({\cal A}_{\eta}\times W)/\mbox{Mod}_{i+1}^{G}({\cal
A}_{\eta}\times W))$. We also want to construct
$\mathfrak{K}^{G}$-counterparts $\mathfrak{M}_{i}^{G}({\cal A}_{\eta}\times
W;\mathfrak{o})$ to the spectra $K^{G}(\mbox{Mod}^{G}_{i}({\cal
A}_{\eta}\times W;\mathfrak{o}))$ where $\mathfrak{o}$ is any orbit under the
action of $G$ on the collection of subsets of cardinality $i$ of
$\\{1,\ldots,n\\}$. In order to do this, we proceed as above by letting
$Q_{i}(\mathfrak{o})$ denote the quotient abelian category
$\mbox{Mod}^{G}({\cal A}_{\eta}\times W)/\mbox{Mod}^{G}_{i}({\cal
A}_{\eta}\times W;\mathfrak{o})$, and defining $\mathfrak{M}_{i}^{G}({\cal
A}_{\eta}\times W;\mathfrak{o})$ to be the homotopy fiber of the map
$\mathfrak{K}^{G}(k)\rightarrow\mathfrak{K}^{G}(Q_{i}(\mathfrak{o}))$. We may
apply the $K^{G}$ case to conclude that there is a decomposition
$\mathfrak{M}_{i}^{G}({\cal A}_{\eta}\times W)/\mathfrak{M}_{i+1}^{G}({\cal
A}_{\eta}\times W)\cong\bigvee_{\mathfrak{o}}\mathfrak{M}_{i}^{G}({\cal
A}_{\eta}\times W;\mathfrak{o})/\mathfrak{M}^{G}_{i+1}({\cal A}_{\eta}\times
W)$ (8–16)
We now have the following Proposition concerning these modules.
###### Proposition 8.1
There are natural maps $\mathfrak{K}^{G}({\cal A}_{\eta}\times W)\rightarrow
Q_{\mbox{dim}(\eta)+1}\rightarrow Q_{1}\rightarrow\mathfrak{K}^{G}({\cal
T}_{\eta}\times W)$.
1. 1.
The map $\mathfrak{K}^{G}({\cal A}_{\eta}\times W)\rightarrow
Q_{\mbox{dim}(\eta)+1}$ is an equivalence of spectra, so we have
$\mathfrak{M}^{G}_{\dim(\eta)+1}({\cal A}_{\eta}\times W)\simeq*$.
2. 2.
The map $Q_{1}\rightarrow\mathfrak{K}^{G}({\cal T}_{\eta}\times W)$ is an
equivalence of spectra, so $\mathfrak{M}_{1}^{G}({\cal A}_{\eta}\times W)$ is
equivalent to the homotopy fiber of the map $\mathfrak{K}^{G}({\cal
A}_{\eta}\times W)\rightarrow\mathfrak{K}^{G}({\cal T}_{\eta}\times W)$.
Proof: For item (1), we note that the category
$\mbox{Mod}^{G}_{\mbox{dim}(\eta)+1}({\cal A}_{\eta}\times W)$ consists only
of the zero module, from which (1) follows. For (2), $Q_{0}$ consists of all
modules on ${\cal A}_{\eta}\times W$ for which the support consists of
coordinate subspaces, including the zero subspace. The localization is
equivalent to the category of free modules over ${\cal T}_{\eta}\times W$.
$\square$
###### Proposition 8.2
In order to prove that the map
$\mathbb{H}_{l}\cong\mathfrak{K}^{G}({\cal A}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l}\xrightarrow{\theta\wedge
id}\mathfrak{K}^{G}({\cal T}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l}$
is an equivalence of spectra, it will suffice to prove that
$\mathfrak{M}^{G}_{i}({\cal A}_{\eta}\times
W;\mathfrak{o})/\mathfrak{M}_{i+1}^{G}({\cal A}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l}\cong*$
for all $i=1,\ldots,n$ and all orbits $\mathfrak{o}$ of the action of $G$ on
the subsets of cardinality $i$ of $\\{1,\ldots,n\\}$.
Proof: Since by Proposition 8.1, $\mathfrak{M}^{G}_{1}({\cal A}_{\eta}\times
W)$ is equivalent to the homotopy fiber of the map
$\mathfrak{K}^{G}({\cal A}_{\eta}\times W)\rightarrow\mathfrak{K}^{G}({\cal
T}_{\eta}\times W)$
it will clearly suffice to prove that
$\mathfrak{M}^{G}_{1}({\cal A}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l}\simeq*$
But we have a diagram of module spectra
$*\simeq\mathfrak{M}^{G}_{dim(\eta)+1}({\cal A}_{\eta}\times
W)\rightarrow\mathfrak{M}^{G}_{dim(\eta)}({\cal A}_{\eta}\times
W)\rightarrow\cdots\rightarrow\mathfrak{M}^{G}_{2}({\cal A}_{\eta}\times
W)\rightarrow\mathfrak{M}^{G}_{1}({\cal A}_{\eta}\times W)$
from which it is clear that it suffices to prove that
$\mathfrak{M}^{G}_{i}({\cal A}_{\eta}\times W)/\mathfrak{M}^{G}_{i+1}({\cal
A}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l}\simeq*$
for $i=1,2,\ldots,dim(\eta)$. The result now follows easily from the
decomposition (8–16) above. $\square$
###### Corollary 8.1
To prove that
$\mathbb{H}_{l}\cong\mathfrak{K}^{G}({\cal A}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l}\xrightarrow{\theta\wedge
id}\mathfrak{K}^{G}({\cal T}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l}$
is an equivalence, it will suffice to prove that the groups
$Tor^{R[G]}_{j}(\pi_{k}(\mathfrak{M}^{G}_{i}({\cal A}_{\eta}\times
W;\mathfrak{o})/\mathfrak{M}_{i+1}^{G}({\cal A}_{\eta}\times
W)),\mathbb{F}_{l})$
vanish for all $i=1,\ldots,n$, and all $j,k$, and $\mathfrak{o}$.
Proof: Follows immediately from Proposition 2.21. $\square$
Let $\eta$ denote an $n$-dimensional $l$-BAR of $G$, and consider an orbit
$\mathfrak{o}$ under the action of $G$ on the subsets of order $i$ of
$\\{1,\ldots,n\\}$. Let $H$ be the stabilizer of one orbit representative in
$\mathfrak{o}$. We now have a commutative diagram
$\begin{diagram}$ (8–17)
and corresponding diagram of representation rings
$\begin{diagram}$ (8–18)
where the quotients $Q_{G}$ and $Q_{H}$ are the images of the composites
$G\stackrel{{\scriptstyle
k}}{{\longrightarrow}}\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n}\rightarrow\Sigma_{n}$
and
$H\stackrel{{\scriptstyle
j}}{{\longrightarrow}}\Sigma_{i}\times\Sigma_{n-i}\ltimes\mathbb{Z}_{l}^{n}\rightarrow\Sigma_{i}\times\Sigma_{n-i}$
The right hand horizontal maps are isomorphisms since the groups $Q_{G}$ and
$Q_{H}$ are finite $l$-groups, by Proposition 6.3. We draw the following
consequence of Proposition 7.2.
###### Proposition 8.3
The $R[G]$ -module $\pi_{k}(\mathfrak{M}^{G}_{i}({\cal A}_{\eta}\times
W;\mathfrak{o})/\mathfrak{M}_{i+1}^{G}({\cal A}_{\eta}\times W))$ is of the
form
$M\stackrel{{\scriptstyle\times
r_{0}}}{{\longrightarrow}}M\stackrel{{\scriptstyle\times
r_{1}}}{{\longrightarrow}}M\stackrel{{\scriptstyle\times
r_{2}}}{{\longrightarrow}}M\stackrel{{\scriptstyle\times
r_{3}}}{{\longrightarrow}}\cdots$
where
1. 1.
$M$ is a module over $R[H]$.
2. 2.
Each $r_{j}$ is a representation of dimension divisible by $l$
3. 3.
Each $r_{j}$ is the pullback along $j$ of a representation $\rho_{i}$ of
$\Sigma_{i}\times\Sigma_{n-i}\ltimes\mathbb{Z}_{l}^{n}$.
Proof: The first follows from Part (1) of Corollary 7.2, with
$M=K^{G_{j}}_{i}(\overline{\cal C}_{\eta})$. The second follows because the
construction $T_{l}$ used in Proposition 7.2 always yields dimension a
multiple of $l$. The last statement follows since the representations
$\sigma_{s}$ are pulled back along a homomorphism to
$\Sigma_{i}\ltimes\mathbb{Z}_{l}^{i}$, which is a subgroup of
$\Sigma_{i}\times\Sigma_{n-i}\ltimes\mathbb{Z}_{l}^{n}$. $\square$
###### Corollary 8.2
There exists an $M$ so that for all $i$, the product $r_{i}\cdot
r_{i+1}\cdot\cdots\cdot r_{i+M-1}\cdot r_{i+M}$ lies in the submodule
$(l)+I_{G}\cdot R[H]$.
Proof: We note that each $\rho_{i}$ has dimension divisible by $l$, and it
readily follows that
$\rho_{i}\in(l)+I_{\Sigma_{i}\times\Sigma_{n-i}\ltimes\mathbb{Z}_{l}^{n}}$
It follows from Corollary 6.4 that there is an $M$ so that
$\rho_{i}\cdot\rho_{i+1}\cdot\cdots\cdot\rho_{i+M-1}\cdot\rho_{i+M}\in((l)+I_{\Sigma_{n}\ltimes\mathbb{Z}_{l}^{n}}+((l)+I_{\Sigma_{i}\times\Sigma_{n-i}})^{\infty})R[\Sigma_{i}\times\Sigma_{n-i}\ltimes\mathbb{Z}_{l}^{n}]$
Referring to the diagram 8–18 above, we have
$j^{*}(\rho_{i}\cdot\rho_{i+1}\cdot\cdots\cdot\rho_{i+M-1}\cdot\rho_{i+M})=r_{i}\cdot
r_{i+1}\cdot\cdots\cdot r_{i+M-1}\cdot r_{i+M}$
and therefore that
$r_{i}\cdot r_{i+1}\cdot\cdots r_{i+M-1}\cdot
r_{i+M}\in((l)+I_{G}+((l)+I_{Q_{H}})^{\infty})R[H]$
But by Proposition 6.3, Part (5), the ideal $((l)+I_{Q_{H}})^{\infty}$ is
equal to the zero ideal, since $Q_{H}$ an $l$-group. It now follows that
$r_{i}\cdot r_{i+1}\cdot\cdots\cdot r_{i+M-1}\cdot r_{i+M}\in(l)+I_{G}R[H]$.
$\square$
###### Corollary 8.3
The groups $Tor^{R[G]}_{t}(\pi_{k}(\mathfrak{M}^{G}_{i}({\cal A}_{\eta}\times
W;\mathfrak{o})/\mathfrak{M}_{i+1}^{G}({\cal A}_{\eta}\times
W)),\mathbb{F}_{l})$ vanish for all $t$.
Proof: Multiplication by an element of $(l)+I_{G}$ annihilates
$Tor^{R[G]}_{i}(M,\mathbb{F}_{l})$. $\square$
This corollary now yields the following.
###### Proposition 8.4
The groups $\pi_{*}(\mathfrak{M}^{G}_{i}({\cal A}_{\eta}\times
W;\mathfrak{o})/\mathfrak{M}_{i+1}^{G}({\cal A}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{\mathfrak{K}^{G}(k)}}\mathbb{H}_{l})$ are
all trivial.
Proof: This is an immediate consequence of Proposition 2.21. $\square$
###### Corollary 8.4
The map $\mathfrak{K}^{G}({\cal A}_{\eta}\times
W)\rightarrow\mathfrak{K}^{G}({\cal T}_{\eta}\times W)$ induces an equivalence
$\mathfrak{K}^{G}(W)\mathrel{\mathop{\wedge}\limits_{K^{G}(k)}}\mathbb{H}_{l}\rightarrow\mathfrak{K}^{G}({\cal
A}_{\eta}\times
W)\mathrel{\mathop{\wedge}\limits_{K^{G}(k)}}\mathbb{H}_{l}\rightarrow\mathfrak{K}^{G}({\cal
T}_{\eta}\times W)\mathrel{\mathop{\wedge}\limits_{K^{G}(k)}}\mathbb{H}_{l}$
Proof: Immediate from Proposition 7.4, and the decomposition 7–11. $\square$
## 9 The main theorem
We are now in position to prove our main theorem.
###### Proposition 9.1
Let $G$ be a totally torsion free profinite $l$-group, and let $W$ be any
Noetherian approximable continuous affine $G$-scheme. Then the natural map of
$\mathfrak{K}^{G}(k)$-module spectra
$\mathfrak{K}^{G}(W)\rightarrow\mathfrak{K}^{G}({\cal E}G\times W)$ induces an
equivalence
$\mathfrak{K}^{G}(W)\mathrel{\mathop{\wedge}\limits_{K^{G}(k)}}\mathbb{H}_{l}\rightarrow\mathfrak{K}^{G}({\cal
E}G\times W)\mathrel{\mathop{\wedge}\limits_{K^{G}(k)}}\mathbb{H}_{l}$
Proof: Since ${\cal E}G$ is an infinite product of continuous affine
$G$-schemes of the form ${\cal T}_{\eta}$, the equivariant
$\mathfrak{K}$-theory of ${\cal E}G\times W$ can be expressed as a colimit of
spectra of the form $\mathfrak{K}^{G}({\cal T}_{\eta}\times W)$. Corollary 8.4
asserts that the map
$\mathfrak{K}^{G}(W)\mathrel{\mathop{\wedge}\limits_{K^{G}(k)}}\mathbb{H}_{l}\rightarrow\mathfrak{K}^{G}({\cal
T}_{\eta}\times W)\mathrel{\mathop{\wedge}\limits_{K^{G}(k)}}\mathbb{H}_{l}$
is an equivalence. The construction
$-\mathrel{\mathop{\wedge}\limits_{K^{G}(k)}}\mathbb{H}_{l}$ commutes with
colimits, and the result therefore follows. $\square$
###### Theorem 9.1
Let $G$ and $W$ be as above, then there are equivalences of spectra
$\mathfrak{K}^{G}(W)^{\wedge}_{\epsilon}\stackrel{{\scriptstyle\sim}}{{\rightarrow}}\mathfrak{K}^{G}({\cal
E}G\times W)^{\wedge}_{\epsilon}\cong\mathfrak{K}({\cal
E}G\mathrel{\mathop{\times}\limits_{G}}W)^{\wedge}_{l}$
well defined up to homotopy, and natural in $W$.
Proof: Follows directly from Proposition 2.21, Corollary 5.3 and Proposition
2.16. $\square$
###### Corollary 9.1
Let $G=G_{F}$ denote the absolute Galois group of a field $F$ containing $k$,
so that $G_{F}$ is a pro-$l$ group and so that $l$ is prime to the
characteristic of $k$. Let $\overline{F}$ denote the algebraic closure of $G$,
equipped with the defining $G_{F}$-action. Then there are natural equivalences
of spectra
$\mathfrak{K}^{G}(k)^{\wedge}_{\epsilon}\stackrel{{\scriptstyle\sim}}{{\rightarrow}}\mathfrak{K}^{G}({\cal
E}G)^{\wedge}_{\epsilon}\cong\mathfrak{K}({\cal E}G/G)^{\wedge}_{l}$
and spectra
$\mathfrak{K}(F)^{\wedge}_{l}\cong\mathfrak{K}^{G}(\overline{F})^{\wedge}_{\epsilon}\stackrel{{\scriptstyle\sim}}{{\rightarrow}}\mathfrak{K}^{G}({\cal
E}G\times Spec(\overline{F}))^{\wedge}_{\epsilon}\cong\mathfrak{K}({\cal
E}G\mathrel{\mathop{\times}\limits_{G}}Spec(\overline{F}))^{\wedge}_{l}$
Proof: We simply apply Theorem 9.1 in the cases $W=Spec(k)$ and
$W=Spec(\overline{F})$. The left hand equivalence in the case
$W=Spec(\overline{F})$ is proved in Corollary 5.4. $\square$
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|
arxiv-papers
| 2013-09-22T21:00:06 |
2024-09-04T02:49:51.322236
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Gunnar Carlsson",
"submitter": "Gunnar Carlsson",
"url": "https://arxiv.org/abs/1309.5649"
}
|
1309.5699
|
# Distribution of Entropy of Bardeen Regular Black Hole with Corrected State
Density
Hai Huang Juhua Chen [email protected] Yongjiu Wang College of Physics
and Information Science, Hunan Normal University, Changsha, Hunan 410081, P.
R. China
###### Abstract
We consider corrections to all orders in the Planck length on the quantum
state density, and calculate the statistical entropy of the scalar field on
the background of the Bardeen regular black hole numerically. We obtain the
distribution of entropy which is inside the horizon of black hole and the
contribution of the vicinity of horizon takes a great part of the whole
entropy.
Keywords: entropy, Bardeen regular black hole
###### pacs:
04.70.Dy, 04.62.+v, 97.60.Lf
In 1970’s, Bekenstein and Hawking found that the black hole(BH) entropy was
proportional to the area of the event horizon by comparing BH physics with
thermodynamics1 .’t Hooft statistically obtained the Bekenstein-Hawking
entropy of a scalar field outside the horizon of the Schwarzschild BH by
brick-wall model2 . The Bekenstein-Hawking entropy was identified with the
statistical-mechanical entropy arising from a thermal bath of quantum fields
propagating outside the event horizon,but there are some drawbacks like
ultraviolet divergence near the horizon. An improved brick-wall method has
been introduced by taking the thin-layer outside the event horizon of a BH as
the integral region, but the artificial cutoffs is still unsolved3 . It’s
commonly believes that there exists a minimal length which originates due to
the quantum fluctuation of the gravitational field motivated by the
generalized uncertainty principle (GUP) 4 ; 5 , the entropy integral of a
radial component in the range near the horizon should be treated for a
convergent entropy6 . The Feynman propagator displays an exponential
ultraviolet cutoff of the form exp($-\lambda p^{2}$), where the parameter
$\lambda$ actually plays the role of a minimal length7 . Quantum gravity
phenomenology has been tackled with effective models based on GUPs and/or
modified dispersion relations containing a minimal length as a natural
ultraviolet cutoff 7.1 ; 7.2 ; 7.3 . At the quantum mechanical level, the
essence of the ultraviolet finiteness of the Feynman propagator can be
captured by a nonlinear relation $p=f(k)$, where $p$ and $k$ are the momentum
and the wave vector of the particle, respectively. The commutator between the
operators $x_{i}$ and $p_{j}$ is generalized to $[x_{i},p_{j}]=i\frac{\partial
p_{i}}{\partial k_{j}}$; moreover, the usual momentum measure in
$3+1$-dimensional spacetime $d^{3}p$ is modified to
$d^{3}p\ \textrm{det}\left|\frac{\partial k_{i}}{\partial
p_{j}}\right|=d^{3}p\prod_{i}\left|\frac{\partial k_{i}}{\partial
p_{j}}\right|,$ (1)
where $\frac{\partial k_{i}}{\partial p_{j}}=\delta_{ij}e^{\lambda
p_{i}^{2}}$8 ; 9 ; 10 , the number of quantum states in a volume element in
phase cell space based on the GUP is
$dg_{+}(\omega)=\frac{d^{3}xd^{3}p}{(2\pi\hbar)^{3}}e^{-\lambda p^{2}},$ (2)
where $p^{2}=p_{i}p^{i}$ is the square of momentum. Many authors calculated
the entropy of BHs to the leading order in the Planck length by using this
newly modified equation of states of density11 ; 12 ; 13 ; 14 . In their
works, the entropy were concentrated on a small vicinity of the BH horizon, we
do not know what’s the contribution of the interior of BH to the entropy.
In this letter, we devote to consider the contribution of the whole interior
of BH by numerical method. We study the entropy of a scalar field on a Bardeen
regular black hole(BRBH) backgrounds. Contrary to our general expectation,
this study of a regular case is not so trivial in contrast to the previous
result of the non-regular case. By using the equation of states of density
motivated by GUP in quantum gravity, we calculate the quantum entropy of a
massive scalar field numerically. We obtain the desired Bekenstein-Hawking
entropy through the whole event horizon and little mass approximation
satisfying the asymptotic property of the wave vector $k$ in the modified
dispersion relation. We take the units as $G=\hbar=c=k_{B}\equiv 1$.
In the context of gravity coupled to some form of matter, regular solutions
are obtained from a prototypical action of the form
$S_{act}=\frac{1}{16\pi}\int d^{4}x\sqrt{-g}(R-\pounds),$ (3)
where $g$ is the determinant of the metric $g_{\mu\nu}$, $R$ is the scalar
curvature, and $\pounds$ represents the Lagrangian of the matter fields. In
the context of a specific nonlinear electrodynamics with Lagrangian given by
$\pounds=(6m/(2|\alpha|\alpha^{2}))(\sqrt{2\alpha^{2}F}/(1+\sqrt{2\alpha^{2}F}))^{5/2}$,
the BRBH line element takes the form
$ds^{2}=-f(r)dt^{2}+\frac{1}{f(r)}dr^{2}+r^{2}(d\theta^{2}+sin^{2}\theta
d\varphi^{2}),$ (4)
note that the function $f(r)$ takes a particularly simple form
$f(r)=1-\frac{2mr^{2}}{(r^{2}+\alpha^{2})^{3/2}},$ (5)
where $(t,r,\theta,\phi)$ are the usual space-time spherical coordinates , $m$
and $\alpha$ stand for the mass and the monopole charge of a self-gravitating
magnetic field of a non-linear electrodynamics source, respectively(seeing
Fig.1). Depending on the value of $\alpha$, the above regular solution may
present two distinct, one,or no horizons. It reduces to the Schwarzschild
solution for $\alpha=0$.
Figure 1: The figure illustrates the behavior of the metric function $f(r)$
given in (4) for different values of the parameter $\alpha$.
In this BRBH background, let us consider massless scalar field which satisfies
the Klein-Gordon equation
$\frac{1}{\sqrt{-g}}\partial_{\mu}(\sqrt{-g}g^{\mu\nu}\partial_{\nu}\Phi)=0.$
(6)
Substituting the line element(4) and using the leading order Wenzel-Kramers-
Brillouin(WKB) approximation $\Phi=e^{-i\omega t}e^{iS(r,\theta,\phi)}$,
Eq.(4) becomes
$\frac{\omega^{2}}{f(r)}-f(r)\left(\frac{\partial S}{\partial
r}\right)^{2}-\frac{1}{r^{2}}\left(\frac{\partial
S}{\partial\theta}\right)^{2}-\frac{1}{r^{2}sin^{2}\theta}\left(\frac{\partial
S}{\partial\varphi}\right)^{2}=0.$ (7)
We have
$p_{r}^{2}=\frac{1}{f(r)}\left(\frac{\omega^{2}}{f(r)}-\frac{1}{r^{2}}p_{\theta}^{2}-\frac{1}{r^{2}sin^{2}\theta}p_{\varphi}^{2}\right),$
(8)
where $p_{r}=\frac{\partial S}{\partial r}$, $p_{\theta}=\frac{\partial
S}{\partial\theta}$ and $p_{\varphi}=\frac{\partial S}{\partial\varphi}$, we
also obtain the square module momentum
$p^{2}=p_{i}p^{i}=g^{rr}p_{r}^{2}+g^{\theta\theta}p_{\theta}^{2}+g^{\varphi\varphi}p_{\varphi}^{2}=\frac{\omega^{2}}{f(r)}.$
(9)
From Eq.(2), the number of quantum states with the energy is less than
$\omega$ is given by
$\displaystyle g_{+}(\omega)$ $\displaystyle=$
$\displaystyle\frac{1}{(2\pi)^{3}}\int e^{-\lambda p^{2}}drd\theta d\varphi
dp_{r}dp_{\theta}dp_{\varphi}$ (10) $\displaystyle=$
$\displaystyle\frac{2}{(2\pi)^{3}}\int e^{-\lambda p^{2}}drd\theta
d\varphi\int\frac{1}{\sqrt{f(r)}}\left(\frac{\omega^{2}}{f(r)}-\frac{1}{r^{2}}p_{\theta}^{2}-\frac{1}{r^{2}sin^{2}\theta}p_{\varphi}^{2}\right)^{\frac{1}{2}}dp_{\theta}dp_{\varphi}$
$\displaystyle=$
$\displaystyle\frac{2}{3\pi}\int\frac{r^{2}\omega^{3}}{f^{2}(r)}e^{-\lambda
p^{2}}dr.$
By using Eq.(10), the free energy can be rewritten as
$\displaystyle F_{+}$ $\displaystyle\approx$ $\displaystyle\frac{1}{\beta}\int
dg_{+}(\omega)\ln(1-e^{-\beta\omega})$ (11) $\displaystyle=$
$\displaystyle-\int\frac{g_{+}(\omega)d\omega}{e^{\beta\omega}-1}$
$\displaystyle=$
$\displaystyle-\frac{2}{3\pi}\int_{0}^{r_{H}}e^{-\frac{\lambda\omega^{2}}{f(r)}}\frac{r^{2}}{f^{2}(r)}dr\int_{0}^{\infty}\frac{\omega^{3}d\omega}{e^{\beta\omega}-1}$
Here, we have taken the continuum limit of quantum numbers and integrated it.
Furthermore, the contribution from the vacuum surrounding the system is
ignored, we are not only concerned about the contribution from the just
vicinity near the horizon, but also the whole region within the horizon
$(0\rightarrow r_{H})$.
We can obtain the entropy as follows:
$\displaystyle S_{+}$ $\displaystyle=$ $\displaystyle\beta^{2}\frac{\partial
F_{+}}{\partial\beta}$ (12) $\displaystyle=$
$\displaystyle\frac{2\beta^{2}}{3\pi}\int_{0}^{r_{H}}e^{-\frac{\lambda\omega^{2}}{f^{2}(r)}}\frac{r^{2}}{f^{2}(r)}dr\int_{0}^{\infty}\frac{e^{\beta\omega}\omega^{4}d\omega}{(e^{\beta\omega}-1)^{2}}$
Figure 2: The figure illustrates the behavior of entropy of massless scalar
wave from BRBH with $\lambda$ from $-0.1$ to $0$ at inverse Hawking
temperature $\beta=25.228$, $\alpha=0.1$, $k_{B}=1$. The bottom picture shows
the entropy together at different Hawking temperature for $\lambda=-0.01$
,$-0.03$, $-0.05$, respectively.
Figure 3: The figure illustrates the behavior of entropy of massless scalar
wave from BRBH with $\lambda$ from $-0.03$ to $0$ at inverse Hawking
temperature $\beta=25.530$, $\alpha=0.2$, $k_{B}=1$. The bottom picture shows
the entropy together at different Hawking temperature for $\lambda=-0.01$
,$-0.03$, $-0.05$, respectively.
Figure 4: The figure illustrates the behavior of entropy of massless scalar
wave from BRBH with $\lambda$ from $-0.1$ to $0$ at inverse Hawking
temperature $\beta=25.133$, $\alpha=0$, $k_{B}=1$. The bottom picture shows
the entropy together at different Hawking temperature for $\lambda=-0.01$
,$-0.03$, $-0.05$, respectively.
Figs.2-4 illustrate the behavior of the BRBH entropy with respect to the
monopole charge $\alpha=0,0.1,0.2$ and parameter $\lambda$, inverse
temperature. We assume $\lambda=\lambda l_{p}^{2}$, where $l_{p}$ is the
Planck length, and in the system of the Planck units $l_{p}=1$. At the same
temperature, the larger the value of $\lambda$, the larger the corresponding
value of entropy is. When $\lambda$ tend to zero, the entropy increase
rapidly. The value of $\lambda$ corresponding to actual entropy is
approximately -0.02. When we used fixed value of $\lambda$, we find the
parameter $\lambda$ makes the entropy lower.
In order to inspect the contribution of different parts of radius inside the
external horizon, we maintain the entropy as $A/4$, and the radius is divided
into 100 parts and integral. In Fig.5, we plot the integral of different parts
for $\alpha=0$, we find the integral increases rapidly when the interval tends
to external horizon. When the interval in the vicinity of zero, it makes
little contribution to the entropy. We show the value of integral for
different $\alpha$ in Table 1, we can see the last parts of integral
constitute the main part of entropy, $\lambda$ approximately equivalent for
different values of temperature. In Fig.6, we show the integral in logarithm
forms to inspect more clearly. When $\alpha=0$(Schwarzschild case), there is
no inner horizon, the entropy increases as integral parts tend to the external
horizon, when $\alpha\neq 0$, there are two extremum. It’s obvious that the
small vicinity of horizon radius show a great contribution than the parts
around horizon radius and the contribution of the vicinity of external horizon
almost takes the whole parts of the entropy. By comparing these curves, we can
see that, the monopole charge $\alpha$ makes the entropy weaker and its
distribution width narrower between two horizon. The inner horizon entropy
increases while the external horizon entropy decrease. Here we adopted the
proper value of $\lambda$, and no cut-off is adopted.
In Fig.7, we try to discuss the situation outside horizon, it shows the
contribution of vacuum outside the horizon. We plot the integral of entropy
from external horizon $r=r_{H}$ to radius $r=100r_{H}$ for
$\alpha=0$(Schwarzschild case),$\lambda=0.0175372$. When $r<10r_{H}$, the
entropy maintains basically stable, when $r>20r_{H}$, the entropy increase as
exponential form. This property indicates that if we make a cut-off between
$r_{H}$ and $r_{H}+\varepsilon$($\varepsilon$ is a small quantity), we can get
the Bekenstein-Hawking entropy. This agrees with the traditional brick-wall
model12 and thin-layer model13 . In Fig.8, if we do not maintain $\lambda$ as
a constant, the integral of entropy from $r_{H}$ to different upper limit can
be adjusted to get the same Bekenstein-Hawking entropy.
Figure 5: The integral of entropy of first 90 parts($0\rightarrow\frac{90}{100}r_{H}$) when $r_{H}$ is divided into 100 parts. The last 10 parts is showed in the table below. Table 1: Entropy Integral between Different Parts of Radius. | $\alpha=0$ | $\alpha=0.1$ | $\alpha=0.2$ | $\alpha=0.3$
---|---|---|---|---
$\lambda$ | -0.0175372 | -0.0175376 | -0.0174500 | -0.0174500
$\beta$ | 25.1327412 | 25.2282438 | 25.5307184 | 26.0959399
$S_{+}$ | 12.5664 | 12.4719 | 12.1855 | 11.6975
91 | 0.00205671 | 0.00204622 | 0.00201396 | 0.00195699
92 | 0.00268403 | 0.00266955 | 0.00262509 | 0.00254671
93 | 0.00360017 | 0.00357973 | 0.00351704 | 0.00340665
94 | 0.00500378 | 0.00497400 | 0.00488275 | 0.00472224
95 | 0.00729502 | 0.00724965 | 0.00711080 | 0.00686681
96 | 0.01137950 | 0.01130580 | 0.01108040 | 0.01068470
97 | 0.01967990 | 0.01954740 | 0.01914320 | 0.01843340
98 | 0.04065840 | 0.04037460 | 0.03951080 | 0.03799340
99 | 0.12414800 | 0.12325100 | 0.12053600 | 0.11575000
100 | 12.2699000 | 12.1775000 | 11.9590000 | 11.481000
Figure 6: The logarithm of the integral of entropy for different interval
within the external horizon $r_{H}$ for $\alpha=0$(Schwarzschild case),$0.1$,
$0.2$, $0.3$ , respectively. Figure 7: The integral of entropy from external
horizon $r=r_{H}$ to radius $r=100r_{H}$ for $\alpha=0$(Schwarzschild
case),$\lambda=0.0175372$. It shows the contribution of vacuum outside the
horizon. Figure 8: The integral upper limit corresponding to different value
of $\lambda$ when entropy is maintain as $A/4$.
In conclusion, we have investigated the entropy of scalar field on a BRBH
backgrounds using analytic and numerical methods. We found the monopole charge
$\alpha$ makes the entropy weaker, and the distribution width of entropy
between two horizons become narrower. We gave the distribution of the entropy
in the interior of BH, we also found that the contribution of the vicinity of
external horizon almost taken the whole parts of entropy. This calculation did
not need to introduce any cut-off while considered the exterior region of
horizon we need to cut out a small region to get the Bekenstein-Hawking
entropy, this agreed with the previous works list in the introduction.
This project is supported by the National Natural Science Foundation of China
under Grant No.10873004, the State Key Development Program for Basic Research
Program of China under Grant No.2010CB832803 and the Program for Changjiang
Scholars and Innovative Research Team in University, No. IRT0964.
## References
* (1) J. D. Bekenstein, Lett. Nuovo Cimento 4 737(1972); J. D. Bekenstein, Phys. Rev. D 7 2333(1973); J. D. Bekenstein, Phys. Rev.D 9 3292(1974); S. W. Hawking, Commun. Math. Phys. 43 199(1975).
* (2) G. ’t Hooft, Nucl. Phys. B256, 727 (1985).
* (3) X. Li and Z. Zhao, Mod. Phys. Letts. A 15 1739(2000); X. Li and Z. Zhao, Int. J. Theor. Phys. 40 90(2001); T. P. Song, C. X. Hou and W. L. Shi, Acta Phys. Sin. 51 06(2002).
* (4) D. J. Gross and P. F. Menda Nucl. Phys. B 303 407(1988).
* (5) Y. W. Kim and Y. J. Park, Phys. Rev. D 77 067501(2008).
* (6) M. Yoon, J. Ha and W. Kim, Phys. Rev. D 76 047501(2007).
* (7) K. Nouicer, Phys. Lett. B 646, 63(2007).
* (8) G. Amelino-Camelia, M. Arzano, Y. Ling, and G.Mandanici, Classical Quantum Gravity 23, 2585 (2006)
* (9) S. Hossenfelder, Phys. Rev. D 73, 105013 (2006); S. Hossenfelder,Class. Quantum Grav. 23, 1815 (2006).
* (10) M. Fontanini, E. Spallucci, and T. Padmanabhan, Phys.Lett. B 633, 627 (2006).
* (11) L. N. Chang, D. Minic, N. Okamura, and T. Takeuchi,Phys. Rev. D 65, 125028(2002).
* (12) Y.W. Kim, Phys. Rev. D 77, 067501(2008).
* (13) X. Li, Phys. Lett. B 540, 9(2002).
* (14) X. Sun and W. Liu, Mod. Phys. Lett. A 19 677(2004).
* (15) R. Zhao, Y. Q. Wu and L. C. Zhang Class. Quantum Grav. 20 4885(2003) C. Liu, X. Li and Z. Zhao, Gen. Rel. Grav. 36 1135(2004); C. Z. Liu, Int. J. Theor. Phys. 44 567(2005).
* (16) W. Kim, Y. W. Kim and Y. J. Park Phys. Rev. D 74 104001(2006); W. Kim, Y. W. Kim and Y. J. Park, Phys. Rev. D 75 127501(2007)
* (17) W. B. Liu, Chin. Phys. Lett. 20 440(2003)
|
arxiv-papers
| 2013-09-23T06:08:21 |
2024-09-04T02:49:51.346276
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Hai Huang, Juhua Chen, Yongjiu Wang",
"submitter": "Juhua Chen",
"url": "https://arxiv.org/abs/1309.5699"
}
|
1309.5717
|
# Ligand-induced stabilization of the aptamer terminal helix in the add
adenine riboswitch
Francesco Di Palma 1, Francesco Colizzi 1 and Giovanni Bussi 1∗
(SISSA - Scuola Internazionale Superiore di Studi Avanzati, via Bonomea 265,
34136 Trieste, Italy)
Stabilization of the P1 stem in add riboswitch
keywords: P1 stem, RNA aptamer, adenine riboswitch, molecular dynamics
simulation, free energy calculation.
11footnotetext: To whom correspondence should be addressed. Email:
[email protected]; Tel: +390403787407
ABSTRACT
Riboswitches are structured mRNA elements that modulate gene expression. They
undergo conformational changes triggered by highly-specific interactions with
sensed metabolites. Among the structural rearrangements engaged by
riboswitches, the forming and melting of the aptamer terminal helix, the so-
called P1 stem, is essential for genetic control. The structural mechanisms by
which this conformational change is modulated upon ligand binding mostly
remain to be elucidated. Here we used pulling molecular dynamics simulations
to study the thermodynamics of the P1 stem in the add adenine riboswitch. The
P1 ligand-dependent stabilization was quantified in terms of free energy and
compared with thermodynamic data. This comparison suggests a model for the
aptamer folding in which direct P1-ligand interactions play a minor role on
the conformational switch when compared with those related to the ligand-
induced aptamer preorganization.
## INTRODUCTION
Riboswitches are ligand-responsive regulatory elements located in untranslated
regions of messenger RNAs (Serganov and Nudler 2013). They change their
conformation in response to specific metabolite binding (Roth and Breaker
2009; Edwards and Batey 2010; Serganov and Patel 2012) and they have been
proposed as modern descendants of an ancient sensory and regulatory system in
the RNA world (Breaker 2012). Many pathogenic bacteria use riboswitches to
control essential metabolic pathways and they are currently regarded as
promising antibacterial drug targets (Blount and Breaker 2006; Mulhbacher et
al. 2010; Deigan and Ferré-D’Amaré 2011). Riboswitches consist of an aptamer
domain that binds the effector ligand and an expression platform that
transduces the ligand-induced conformational “switch” into a modulation of
gene expression (Barrick and Breaker 2007; Garst and Batey 2009). Among more
than 20 natural aptamer classes (Breaker 2012), purine-sensing riboswitches
have the peculiarity to recognize the targeted purine by utilizing a conserved
pyrimidine (Kim and Breaker 2008; Batey 2012). One of the most characterized
members of this class is the adenine sensing riboswitch (A-riboswitch) _cis_
-regulating the _add_ gene in _Vibrio vulnificus_ (Mandal and Breaker 2003).
The ligand-bound structure of its aptamer is a three-way junction composed of
three stems (P1, P2, P3) with the ligand completely encapsulated into the
structure (Fig. 1) (Serganov et al. 2004; Mandal and Breaker 2004). The
specificity for adenine is ensured by canonical Watson-Crick (WC) base pairing
established between a uracil in conserved position and the ligand (Noeske et
al. 2005; Gilbert et al. 2006).
The A-riboswitch acts as a translational regulator (Serganov et al. 2004;
Lemay et al. 2011). In the absence of adenine the ribosome binding site and
the initiation codon, which are portions of the expression platform, are
sequestered by pairing with a portion of the aptamer (OFF-state, Fig. 2B). The
presence of adenine stabilizes an aptamer conformation in which the terminal
P1 helix is well structured and both the regulatory sequences are available
for ribosomal binding thus enabling mRNA translation (ON-state, Fig. 2A)
(Rieder et al. 2007; Lee et al. 2010; Leipply and Draper 2011). The structural
mechanism regulating the switch between the ON- and the OFF-state upon ligand
binding mostly remains to be elucidated. The P1 stem is formed in the ON-state
and it is disrupted in the OFF-state (Mandal and Breaker 2003; Serganov et al.
2004).
It has been proposed that P1 is stabilized by the ligand (Batey et al. 2004)
and that this could be a common feature in many riboswitch classes (Montange
and Batey 2006). The role of ligand binding in the structural organization of
the aptamer has been investigated with single-molecule spectroscopy providing
an insightful overview on the folding dynamics (Neupane et al. 2011), yet
lacking the critical atomistic details needed for an accurate structural
characterization of the process as extensively discussed by Lin et al. (2012).
Although _in silico_ techniques have been used to investigate the ligand role
(Lin and Thirumalai 2008; Sharma et al. 2009; Priyakumar and MacKerell 2010;
Gong et al. 2011; Allnér et al. 2013), a quantitative estimation of the
energetic contributions associated to ligand binding, in particular regarding
the role of direct P1-ligand interactions, has not yet been provided. In this
context, state-of-the-art free-energy methods combined with atomistic
simulations can bridge the gap providing an unparalleled perspective on the
mechanism and dynamics of the biomolecular process of interest (Dellago and
Bolhuis 2009). In this work we used steered molecular dynamics (SMD)
(Grubmuller et al. 1996; Sotomayor and Schulten 2007) simulations to study the
thermodynamics of the P1 stem formation in the presence and in the absence of
the cognate ligand. We enforced the breaking of the P1 stem base pairs (bp)
and then using a recently developed reweighting scheme (Colizzi and Bussi
2012) we quantitatively estimated the ligand-induced stabilization of the
helix. The A9-U63 bp which directly stacks with adenine was used as a proxy
for the P1 stability. Our non-equilibrium simulations provide measurements of
the stability of the A9-U63 bp and quantify the direct ligand-dependent
stabilization of the pairing. In the following our results are presented and
compared with melting and single-molecule experiments. A structural model for
the conformational switch emerging from the combination of our results and
previous experimental data is also discussed.
## RESULTS
We carried out the simulations of the aptamer domain of the _add_ A-riboswitch
in different forms, namely the entire aptamer (PDB id 1Y26) (Serganov et al.
2004) has been simulated in the presence (Holo) and in the absence (Apo) of
the cognate ligand, the adenine; additionally, to better estimate the ligand-
induced stabilization, we also simulated a truncated aptamer
($\triangle$1-8/64-71), both in the Apo and Holo forms. Long unbiased
molecular dynamics (MD) for all the four systems were performed to test the
stability of the aptamer in different conditions. In the truncated systems the
terminal bp was restrained in its initial configuration to mimic the presence
of the rest of the stem. Furthermore, the full-length systems were pulled from
the terminal bases to disrupt the entire P1 stem (Fig. 3) thus allowing its
different stability between in the Holo and in the Apo forms to be
qualitatively inferred. At last, to quantify this difference, SMD simulations
of both the $\triangle$1-8/64-71 systems were done enforcing the breaking of
the A9-U63 bp that directly stacks with the ligand (Fig. 4).
The stability of both the Apo and Holo systems was evaluated monitoring the
root mean square deviation (RMSD) from the native structure along 48 ns MD
runs (Fig. 5A-B). Ligand removal (see Materials and Methods for details) did
not affect the overall stability of the Apo aptamer in this time-scale, and
secondary and tertiary structures were substantially unchanged.
The analysis of the trajectories obtained by pulling the P1 stem showed that
the secondary and tertiary structure elements of the rest of the aptamer were
not affected by the opening of the helix (data not shown). Focusing our
attention on the P1 stem we observed that in the Apo form the A9-U63 bp (Fig.
3) was broken when the distance between the centers of mass of the terminal
bases reached a value of $\approx$9.8 nm. Differently, in presence of the
ligand a longer pulling was needed and the rupture only happened at a distance
of $\approx$11.5 nm (Fig. 6). This behavior is compatible with the picture in
which the ligand stabilizes the P1 stem (Montange and Batey 2006). It was
however difficult to extract quantitative information on the ligand-P1
interaction from these simulations because the rupture is a stochastic event
and extensive sampling would be required. Moreover, as pointed out in a recent
paper (Lin et al. 2012), the end-to-end distance could be a non-optimal CV for
pulling experiments or simulations since local bp formation plays an important
role in global stem folding.
The quantitative analysis of the P1-ligand interaction was better obtained
from the simulation of both the $\triangle$1-8/64-71 systems. We verified
that, when the P1 stem is replaced with the A9-U63 bp restrained to be in
canonical WC pairing, the aptamer remains stable (Fig. 5C-D). Remarkably,
fluorescence experiments have shown that the aptamer can also fold and bind
adenine when large fractions of the P1 stem are removed (Lemay and Lafontaine
2007). This validates the possibility of using the two structures,
$\triangle$1-8/64-71 Holo and $\triangle$1-8/64-71 Apo, to investigate the
direct P1 stabilization given by the adenine. In the following we focus on the
SMD simulations performed on these truncated forms. Typical initial and final
conformations from the SMD are shown in Fig. 4.
### Analysis of work profiles
The unbinding event of the A9-U63 bp is described as a function of the value
of the steered RMSD in Fig. 7. The initial value corresponds to the
configuration with the WC pairing formed, whereas at the final value (0.35 nm)
the pairing is completely broken. Even if the ensembles of work profiles for
the two forms are broadly spread and superposed, the free-energy profiles
computed using the Jarzynski equality (Jarzynski 1997 _a_) as the exponential
average of the two sets of data are clearly distinguishable (Fig. 7A).
Qualitatively it is worth highlighting that during the breaking of the A9-U63
bp the Apo form (red line) is always lower in free energy than the Holo form
(blue line). It follows that the breaking of the monitored bp in the Apo form
was unambiguously more probable than in the Holo one (Fig. 7A). However, such
an approach was still a long way off from quantitatively accounting for the
energetic stabilization of the A9-U63 bp related to the presence of adenine in
the binding site. Within this framework there was no way to automatically
detect when the nucleobases reached the unbound configuration. It was thus
difficult to avoid systematic errors in the comparison of the two systems.
Furthermore, few low-work realizations occurred during the unpairing in the
presence of adenine. In these low-work realizations the number of hydrogen
bonds was non-zero at large RMSD values and the structural analysis of the
trajectories revealed the transient formation of a cis-sugar edge pair (data
not shown) (Leontis et al. 2002). Due to the exponential nature of the
Jarzynski average, these low-work realizations dominated the free energy
profile for the Holo form further compromising the possibility of a
quantitative comparison with the Apo form.
### Energetics of hydrogen bond breaking
We thus analyzed the trajectories in terms of number of hydrogen bonds formed
between A9 and U63, a discrete variable that more strictly reported on the
breaking of the pairing. In this metrics, the bound (1 or 2 hydrogen bonds)
and unbound (0 hydrogen bond) ensembles could be clearly and unambiguously
identified thus allowing a quantitative comparison between the Apo and the
Holo system. Additionally, the configurations from the outlier trajectories
could be assigned properly to one or the other ensemble in spite of their
atypical RMSD value.
The differences in free energy ($\Delta F$) between the ensembles, with and
without hydrogen bonds, was computed using a reweighting scheme (Colizzi and
Bussi 2012). The values and the associated standard errors were estimated for
both systems. For the Apo form $\Delta F=-2.5\pm 1.4\ kJ/mol$ suggesting that
the bp could spontaneously break in the absence of adenine. For the Holo form
$\Delta F=1.9\pm 1.7\ kJ/mol$ implying that the presence of the ligand and its
pairing with U63 stabilized the stacked bp. The $\Delta\Delta F$ between the
two forms is equal to $-4.4\pm 2\ kJ/mol$. This value quantifies the
thermodynamic stabilization to the formation of the base pair which directly
interacts with adenine in the P1 stem.
## DISCUSSION AND CONCLUSIONS
Our simulations at atomistic detail provide for the first time the free-energy
contribution of ligand stacking to the formation of the P1 stem in a
riboswitch. In particular, the presented _in silico_ approach allows the
energetics involved in the aptamer stabilization upon ligand binding to be
dissected in detail. Below we compare our results with single-molecule
manipulation, both _in vitro_ and _in silico_ , and thermodynamic data from
dsRNA melting experiments. We also provide a model for ligand-modulated co-
transcriptional folding of the _add_ riboswitch.
### Comparison with related works
Our results are in nice agreement with thermodynamic data based on dsRNA
melting experiments (Mathews et al. 2004; Turner and Mathews 2010). The
comparison between our simulations and those experiments can be
straightforwardly achieved by considering the pairing between U62 and the
sensed adenine as an additional terminal bp of the P1 stem. The direct
stabilization of the P1 stem due to the cognate-ligand binding should be thus
equivalent to that given by adding one further AU bp to the P1 helix. Using
the most recent nearest neighbor energy parameters for the comparison of RNA
secondary structures (Mathews et al. 2004; Turner and Mathews 2010), the free-
energy difference between the sequence of the P1 stem with and without the
additional AU base pair, ${}_{\textrm{5'-CGCUUCAUA\emph{A}-3'}}$
${}^{\textrm{3'-GUGAAGUAU\emph{U}-5'}}$ and ${}_{\textrm{5'-CGCUUCAUA-3'}}$
${}^{\textrm{3'-GUGAAGUAU-5'}}$ , can be computed (Hofacker et al. 1994;
Lorenz et al. 2011) as $\Delta\Delta F=-3.7\ kJ/mol$, consistently with our
results.
Our free-energy estimates complement previously reported investigations in
which the role of the ligand in the folding process of the A-riboswitch has
always been referred to the whole aptamer (Lin and Thirumalai 2008; Neupane et
al. 2011) and never specifically to the P1 stem. Using a one-bead-per-
nucleotide coarse-grained model, the $\Delta\Delta F$ has been computed as
approximately $-15\ kJ/mol$ (Lin and Thirumalai 2008). Notably also this
calculation has been done using a shortened P1 stem, possibly affecting the
$\Delta F$ estimation. Single-molecule force spectroscopy experiments have
been also performed to characterize the folding pathway of the aptamer with an
estimated $\Delta\Delta F\simeq-33\ kJ/mol$ (Neupane et al. 2011). However in
both these works the separated contributions of the P1-ligand stacking, of the
interaction between the ligand and the junctions J1-2, J2-3 and J3-1, and of
the interaction between loops L2-L3 could not be discerned (secondary
structure elements labeled as in Fig. 8).
From the comparison of our data with the above mentioned experimental and
computational works, a twofold modular role for the ligand emerges. On the one
hand, the binding of adenine can contribute to the aptamer preorganization and
it could allow the long-range induction of the tight hydrogen-bonding and
base-stacking networks observed in the native state (Rieder et al. 2007; Lee
et al. 2010). This preorganization would reduce the distance between A9 and
U63, thus increasing the probability of their pairing. A similar mechanism has
been proposed also for the SAM-I riboswitch (Whitford et al. 2009). On the
other hand, adenine binding enhances the P1 formation by direct stacking
interaction, mimicking the extension of the helix by an additional bp.
Notably, the energetic contribution of the direct stacking is smaller than
that involved in the aptamer preorganization. The latter can be estimated as
the difference between the global ligand-induced aptamer stabilization (Lin
and Thirumalai 2008; Neupane et al. 2011) and the stacking contribution
dissected in our work.
### Folding model
Our work provides atomistic details and energetic estimates to the currently
accepted model for the folding of the _add_ riboswitch upon ligand binding
(Rieder et al. 2007; Lee et al. 2010; Leipply and Draper 2011). Altogether,
our data and the related experimental works suggest a folding model as
depicted in Fig. 8. Initially, only the P2 and P3 stems and the corresponding
loops (L2, L3, still not interacting each other) are formed and not fully
stable (Fig. 8A). Then, adenine binding allows for a preorganization of the
aptamer where the three junctions arrange around the ligand (Fig.
8$\textrm{B}_{\textrm{1}}$), stabilizing also the previously formed helices
(Rieder et al. 2007). It has not been established clear if the loop-loop
interaction is formed before or after ligand binding (Leipply and Draper
2011). Thus, an alternative pathway, the junctions and the P1 could acquire a
partially folded conformation also in the absence of adenine (Fig.
8$\textrm{B}_{\textrm{2}}$) (Lee et al. 2010). Finally, the P1 helix becomes
fully structured and stabilized by the ligand (Fig. 8C), to the detriment of
the expression platform (see Fig. 2) (Lemay et al. 2011). This step is
mandatory for translation to be initiated. We quantified the ligand
contribution to the P1 stem formation due to direct interactions to be
approximately -4 kJ/mol.
Our result is compatible with both the folding pathways (Figs.
8$\textrm{B}_{\textrm{1}}$ and 8$\textrm{B}_{\textrm{2}}$) irrespectively of
their relative population and cannot discriminate among them. The relative
probability between the two paths can be modulated by the ligand concentration
and its binding affinity. On the one hand, the intermediate shown in Fig.
8$\textrm{B}_{\textrm{1}}$ could be relevant for ligand-RNA binding in an
early transcriptional context in which the last 9 nucleobases (i.e those
allowing P1 formation) of the aptamer have not yet been synthesized. Indeed,
it has been shown that also an aptamer missing a large portion of the P1 stem
is able to bind adenine (Lemay and Lafontaine 2007). On the other hand, the
intermediate shown in Fig. 8$\textrm{B}_{\textrm{2}}$ could be populated at
low-ligand concentration once the nucleobases allowing P1 formation are
synthesized. Later on, after the synthesis of the expression platform, ligand
binding could shift the thermodynamic equilibrium towards one of two competing
riboswitch conformations (P1 formed and non formed).
### Conclusion
Ligand-induced stabilization of the P1 stem is crucial for A-riboswitch
regulation and function. Here we quantified the direct interaction between
adenine and P1 stem and analyzed it at atomistic detail. Our results suggest a
model for the aptamer folding in which the direct P1-ligand interactions play
a minor role on the conformational switch when compared with those related to
the ligand-induced aptamer preorganization. Because the structural/functional
role of the aptamer terminal helix is a common feature in the “straight
junctional” riboswitches (Serganov and Nudler 2013), we foresee a wider
validity of the model presented herein.
## MATERIALS AND METHODS
### System description and set-up
We simulated the Holo and the Apo form of the A-riboswitch aptamer domain both
composed of 71 nucleotides. The Apo form was generated by adenine removal from
the ligand-bound (Holo) crystal structure (PDB id: 1Y26) (Serganov et al.
2004). This deletion is justified by the fact that the Apo and Holo form have
been shown experimentally to share an overall similar secondary structure
(Lemay et al. 2011). This is at variance with e.g. the _pbuE_ adenine
riboswitch in which the two structures are different. The generation of the
Apo form by simply removing the ligand has been adopted also in a recent work
(Allnér et al. 2013). Molecular dynamics (MD) simulations were performed using
the Amber99 force field (Wang et al. 2000) refined with the _parmbsc0_
corrections (Pérez et al. 2007). From the analysis of the SMD trajectories we
do not expect the refinement on the $\chi$ torsional parameters (Zgarbová et
al. 2011) to affect the results. Adenine was parametrized using the general
Amber force field (gaff) (Wang et al. 2004). Partial atomic charges were
assigned using the restricted electrostatic potential fit method (Bayly et al.
1993) based on an electronic structure calculation at the HF/6-31G* level of
theory performed with Gaussian03 (Frisch et al. 2004). Bond-lengths were
constrained with LINCS (Hess et al. 1997) and the electrostatic interactions
were calculated using the particle-mesh Ewald method (Darden et al. 1993). For
both forms, the following protocol was used to prepare the systems (Table 1
for details) for MD simulations: steepest descent minimization (200 steps)
starting from the X-ray structure. Solvation with $\approx$13000 TIP3P water
molecules (Jorgensen et al. 1983) and NaCl at 0.15 M concentration (plus extra
Na+ counter-ions to neutralize the charges of the systems) in a hexagonal
prism (lattice vectors in nm [(10,0,0), (0,7,0),
(0,$\frac{7}{2}$,$\frac{7\cdot\sqrt{3}}{2}$)]) that was created orienting the
major length of the aptamer along the X axis. Steepest descent minimization
(200 steps) for ions and solvent; the systems were thermalized at 300 K,
initially for 200 ps with frozen solute positions and then for 5 ns in NPT
ensemble (1 atm) with stochastic velocity rescaling (Bussi et al. 2007) and
Berendsen barostat (Berendsen et al. 1984); to maintain the systems oriented
along the largest lattice vector (X) a restraint was imposed with a force
constant of
4$\cdot$103$\nicefrac{\mathrm{(\nicefrac{\mathrm{kJ}}{\mathrm{mol}})}}{\mathrm{nm^{2}}}$
on the Y and Z components of the distance between phosphate atoms of A52 and
G71. Each system was simulated for 48 ns in NVT ensemble to assess the
stability of the aptamer.
### Steered molecular dynamics
To perform SMD simulations inducing the opening of the whole P1 stem the
systems were solvated again with $\approx$39500 water molecules in a larger
rhombic dodecahedral box with distance between periodic images equal to 12 nm,
adding ions to maintain the same ionic strengh (P1-SMD systems in Table 1).
The same protocol described above was applied for the minimization,
thermalization and equilibration of this larger Holo and Apo systems for the
pulling simulations. An incremental separation between the centers of mass of
the terminal nucleotides (C1 and G71) was imposed from an initial value of
1.05 nm to a value sufficient to completely unfold the 9 bp of the P1 helix
(Apo 10.05 nm, Holo 11.75 nm) at a speed of 0.56 nm/ns (see Fig. 3). The
spring constant was set to
3.9$\cdot$104$\nicefrac{\mathrm{(\nicefrac{\mathrm{kJ}}{\mathrm{mol}})}}{\mathrm{nm^{2}}}$.
The first eight bp of the P1 stem (i.e. whole stem from C1 to U8 and from A64
to G71, except for A9-U63 bp) were then cut in both systems creating the
$\triangle$1-8/64-71 Holo and $\triangle$1-8/64-71 Apo structures (Table 1 for
details). Water molecules were allowed to relax filling the space left by the
16 removed bases through an additional 1 ns equilibration in which the
positions of aptamer atoms were frozen followed by 5 ns of unrestrained NPT
simulation. Then the systems were simulated for 48 ns in the NVT ensemble
restraining the terminal bases in the initial state to avoid any spontaneous
flipping. The pairs deletion is not biologically meaningless, because it has
been shown experimentally that a series of aptamer variants with shorter P1
helix are still able to bind the ligand (Lemay and Lafontaine 2007). The
deletion reduced the noise during the pulling allowing to focus the
calculation on the influence of the ligand on the A9-U63 pairing. This bp
rupture was here enforced by pulling on the RMSD between the heavy atoms of A9
and U63 with reference to the crystal structure. This collective variable (CV)
was chosen as it identifies the native conformation (RMSD $\approx$ 0) of the
A9-U63 bp, which is necessary for the initiation of the P1 stem formation. The
steered CV was pulled at constant velocity of 0.175 $\nicefrac{{nm}}{{ns}}$
from 0 to 0.35 nm in 2 ns. This pulling induced the complete opening of the
A9-U63 bp in presence and absence of the ligand (Fig. 4). The spring constant
was set to
3.9$\cdot$104$\nicefrac{\mathrm{(\nicefrac{\mathrm{kJ}}{\mathrm{mol}})}}{\mathrm{nm^{2}}}$.
The starting points were extracted equidistantly (one every 16 ps) from a
8.192 ns run NVT ensemble restraining the RMSD value of those atoms at 0. For
the two system 512 independent SMD simulations were performed, corresponding
to an aggregate time of approximately 1$\mu$s each. Simulations were carried
out with the Gromacs 4.0.7 program package (Hess et al. 2008) combined with
the PLUMED 1.3 plug-in (Bonomi et al. 2009).
### Analysis
The Jarzynski equality (Jarzynski 1997 _a_) was used to estimate the
equilibrium free-energy landscape of the process from the collected work
profiles. The simulations were then analyzed using a recently proposed scheme
(Colizzi and Bussi 2012) which combines an identity by Jarzynski (Jarzynski
1997 _b_) with the weighted-histogram analysis method (Kumar et al. 1992). The
algorithm allows the free-energy profiles to be projected onto any _a
posteriori_ chosen CV. It is well known that free-energy calculations using
Jarzynski-based relationships are difficult to converge. Statistical errors
were thus estimated by the bootstrapping procedure described in Do et al.
(2013) indicating that our results were converged within $\approx k_{B}T$. The
VIENNA RNA package (Hofacker et al. 1994; Lorenz et al. 2011) was used to
compare our results with the thermodynamic data based on dsRNA melting
experiments (Mathews et al. 2004; Turner and Mathews 2010).
## ACKNOWLEDGMENTS
We thank Daniel Lafontaine and Gabriele Varani for reading the manuscript and
providing critical comments. Sandro Bottaro is also acknowledged for carefully
reading the manuscript and suggesting several improvements. We acknowledge the
CINECA award no. HP10B2G6OF, 2012 under the ISCRA initiative for the
availability of high performance computing resources. The research leading to
these results has received funding from the European Research Council under
the European Union’s Seventh Framework Programme (FP/2007-2013) / ERC Grant
Agreement n. 306662, S-RNA-S.
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Table 1: Specifications for the simulated systems System | Total atoms | Water molecules | aptamer atoms | Ions (Na+Cl)
---|---|---|---|---
Holo form | 41628 | 13078 | 2257 | 122 (96+26)
Apo form | 41676 | 13099 | 2257 | 122 (96+26)
Holo (P1-SMD) | 120654 | 39364 | 2257 | 290 (180+110)
Apo (P1-SMD) | 120588 | 39347 | 2257 | 290 (180+110)
Holo $\Delta$1-8/64-71 | 41101 | 13078 | 1746 | 106 (80+26)
Apo $\Delta$1-8/64-71 | 41149 | 13099 | 1746 | 106 (80+26)
Figure 1: Adenine riboswitch aptamer and binding site. A) Secondary structure
elements and B) 3-dimensional structure with bound adenine. The P1 stem is
grey, the other stems and loops are black. C) Cartoon representation of the
binding site; the two dotted lines represent the hydrogen bonds of the WC
pairing between the U62 and the ligand.
Figure 2: Secondary structure representation of the _add_ riboswitch in the
ON (A) and OFF (B) states. The ligand, the initiation codon and the Shine-
Dalgarno sequence are labeled.
Figure 3: Initial (A) and final (B) configuration of the SMD simulation
opening the P1 stem here shown for the Holo form. The backbone of the aptamer
is in blue except for the P1 stem, in light blue. The ligand and the 18 bases
forming the helix are shown. The P1 stem is formed in A and disrupted in B.
Figure 4: Representative structures of the Holo binding pocket at the
beginning (A, RMSD = 0 nm) and at the end (B, RMSD = 0.35 nm) of the SMD. The
portion of the P1 stem removed in our simulations is in light blue. Bases
forming the binding pocket are labeled, ligand is shown in red. A9-U63 pair is
formed in A and disrupted in B.
Figure 5: Root mean square deviation (RMSD) from native structure. A) Holo
form during 48 ns equilibration, computed on the whole aptamer (black) and on
the bases from 9 to 63 (grey). B) Same as A) done on the Apo form (whole
aptamer, black; bases from 9 to 63, grey). The difference between black and
grey profiles in both panels, indicates that the P1 stem is less stable than
the rest of the aptamer. C) $\triangle$1-8/64-71 Holo RMSD along the NVT 48 ns
equilibration. D) Same as C) for the $\triangle$1-8/64-71 Apo form.
Figure 6: Base-pair ruptures during P1 pulling. In the pulling simulations
the 9 bp forming the P1 stem were unpaired. We here monitored the RMSD of each
bp (grey-scale) from their native conformation (A, Holo; B, Apo). A9-U63 bp
(in black) was disrupted (RMSD $\approx$0.5 nm, arrows) later in the Holo form
($\approx$19 ns) than in the Apo one ($\approx$16 ns).
Figure 7: Unbinding process of the A9-U63 bp. (A) Mechanical work performed
as a function of the value of the steered RMSD, or equivalently of time, for
512 simulations for Apo (pink) and Holo form (light-blue). The respective free
energies resulting from the Jarzynski equality are shown in thicker red and
blue lines. The initial free-energy decrease related to the entropy gain
induced by the restraint movement has no consequence on the final result. (B)
Hydrogen bonds occurrence for both the systems Apo (red) and Holo (blue). (C)
Snapshots of the Holo form (ligand in black) with 2, 1 or 0 hydrogen bond
(dotted lines) formed between A9-U63.
Figure 8: Schematic representation of the aptamer secondary structure in its
folding intermediates with and without the ligand. A) Stems P2 and P3, loops
L2 and L3 are folded but not stable. $\textrm{B}_{\textrm{1}})$ The junctions
(J1-2, J2-3, J3-1) arrange around the ligand (ADE) and the inter-loop pairings
occur (L2-L3) stabilizing also the stems. $\textrm{B}_{\textrm{2}})$
Alternative possible intermediate in which all the 3 stems are not completely
and stably folded before ligand binding. The junctions and the tertiary
interaction between the loops are not stable. C) P1 stem is fully folded and
stabilized by the ligand.
|
arxiv-papers
| 2013-09-23T07:27:52 |
2024-09-04T02:49:51.353122
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Francesco Di Palma, Francesco Colizzi, and Giovanni Bussi",
"submitter": "Giovanni Bussi",
"url": "https://arxiv.org/abs/1309.5717"
}
|
1309.5767
|
Hugo Herbelin and Arnaud Spiwack Billy Editor, Bill Editors2Conference title
on which this volume is based on111 10.4230/LIPIcs.xxx.yyy.p
# The Rooster and the Syntactic Bracket111This research has received funding
from the European Research Council under the FP7 grant agreement 278673,
Project MemCAD
Hugo Herbelin Inria Paris-Rocquencourt
Paris, France
[email protected] Arnaud Spiwack Inria Paris-Rocquencourt
Paris, France
[email protected]
###### Abstract.
We propose an extension of pure type systems with an algebraic presentation of
inductive and co-inductive type families with proper indices. This type theory
supports coercions toward from smaller sorts to bigger sorts via explicit type
construction, as well as impredicative sorts. Type families in impredicative
sorts are constructed with a bracketing operation. The necessary restrictions
of pattern-matching from impredicative sorts to types are confined to the
bracketing construct. This type theory gives an alternative presentation to
the calculus of inductive constructions on which the Coq proof assistant is an
implementation.
###### Key words and phrases:
Coq, Calculus of inductive constructions, Impredicativity, Strictly positive
type families, Inductive type families
###### 1991 Mathematics Subject Classification:
F.3.3 Studies of Program Constructs
## 1\. Introduction
In the Coq proof assistant [14] inductive types are treated as toplevel
definitions. If it makes sense from a convenience or an efficiency point of
view, the monolithic nature of the definitions make it hard to describe what
they precisely mean. As a matter of fact, inductive definitions mean different
things depending on the type they are defined in: specifically, some types are
interpreted differently in impredicative sorts like Prop or the impredicative
variant of Set.
In this article, we present a calculus of inductive and co-inductive
constructions where inductive and co-inductive types are presented
algebraically. The algebraic presentation is an extension of a PTS [3] with
inductive and co-inductive type families. Thanks to its modularity, it is
meant to serve as a description which is simpler to expose and more
mathematically amenable than the monolithic scheme which is found in a
practical system such as Coq. For the sake of clarity, the system is given
with a single universe and explicit subtyping, although Coq has an unbounded
cumulative hierarchy of universes and implicit subtyping. Apart from these
technicalities, it is believed that our calculus of algebraic inductive and
co-inductive constructions expresses all the features of the Set-impredicative
Calculus of Inductive Constructions that Coq implements, _e.g._ in its version
8.4 when launched with option -impredicative-set.
This work draws most of its inspiration from Morris & al [9, 10] for the
algebraic presentation of inductive type families in a predicative sort, and
Awodey & Bauer [2] for the treatment of impredicative sorts.
We use examples from Coq to illustrate the algebraic presentation. To
differentiate expressions in Coq from expressions in the algebraic
presentation, the former are typeset in a sans-serif font and the latter in a
roman font.
## 2\. Pure type systems
To model the type system of Coq, we start with the classic presentation of
pure type systems (pts) of Barendregt [3], which we will then extend to model
type families. A pts is characterised by a _single_ syntactic category of
terms which are used both as ${\lambda}$-terms and as types. It has a single
form of typing judgment $\Gamma\vdash u:A$, where $u$ and $A$ are terms, and
$\Gamma$ a context assigning terms to variables. A pts has a set of _sorts_ ,
which we shall denote schematically by the symbol $s$. Every sort is an atomic
term. A pts has a conversion relation $u\equiv v$. Here we diverge from the
presentation of [3] which always uses ${\beta}$-conversion. Coq, on the other
hand, uses ${\beta}$${\eta}$-conversion on the fragment described in this
section. In this article we will take the conversion rule as abstract, not
even requiring it to be decidable. We will only require that it contains all
the reduction rules which are given in the form $u\leadsto v$ (in this
section, we only have ${\beta}$-reduction).
$\Gamma,x{:}A\vdash x:A\lx@proof@logical@and\Gamma\vdash A:sx\mbox{{ is fresh
in }}\Gamma$ $\Gamma,x{:}B\vdash u:A\lx@proof@logical@and\Gamma\vdash
u:A\Gamma\vdash B:sx\mbox{{ is fresh in }}\Gamma$
$\Gamma\vdash\lambda{x}^{A}.\,u:\prod_{x{:}A}B\lx@proof@logical@and\Gamma\vdash\prod_{x{:}A}B:s\Gamma,x{:}A\vdash
u:B$
$\Gamma\vdash{u}\,{v}:{B}{\left[\,x\,{\setminus}\,v\,\right]}\lx@proof@logical@and\Gamma\vdash
u:\prod_{x{:}A}B\Gamma\vdash v:A$
$\Gamma\vdash u:B\lx@proof@logical@and\Gamma\vdash u:A\Gamma\vdash B:sA\equiv
B$
${\left(\lambda{x}^{A}.\,u\right)}\,{v}\leadsto{u}{\left[\,x\,{\setminus}\,v\,\right]}$
Figure 1. Generic rules of pts
The typing rules of a pts comprise of a set of generic rules given in Figure
1, together with a number of rules of the form
$\vdash{s}_{1}:{s}_{2}$
called axioms, and rules of the form
$\Gamma\vdash\prod_{x{:}A}B:{s}_{3}\lx@proof@logical@and\Gamma\vdash
A:{s}_{1}\Gamma,x{:}A\vdash B:{s}_{2}$
of product formation rules. As usual we write $A\rightarrow B$ for
$\prod_{x{:}A}B$ when $x$ does not bind a variable in $B$.
As a starting point of the algebraic presentation, we shall use a pts with two
sorts, $\mathrm{Type}$ and $\Box$, together with the following axiom:
$\Gamma\vdash\mathrm{Type}:\Box$
and the following product formation rule:
$\Gamma\vdash\prod_{x{:}A}B:{\mathrm{max}}\,{{s}_{1}}\,{{s}_{2}}\lx@proof@logical@and\Gamma\vdash
A:{s}_{1}\Gamma,x{:}A\vdash B:{s}_{2}$
where ${\mathrm{max}}\,{s}\,{s}=s$ and
${\mathrm{max}}\,{\mathrm{Type}}\,{\Box}={\mathrm{max}}\,{\Box}\,{\mathrm{Type}}=\Box$
The sorts $\mathrm{Type}$ and $\Box$ are predicative. The sort $\Box$ plays a
technical role in allowing type variable and the formation of type-level
functions; it cannot, however, be referenced in terms. In the following
sections, $\Box$ will also be used to be able to define types by pattern-
matching (_strong elimination_).
To model the entire Coq system, $\mathrm{Type}$ and $\Box$ would be replaced
with a hierarchy of predicative sorts ${\mathrm{Type}}_{i}$, such that
$\Gamma\vdash{\mathrm{Type}}_{i}:{\mathrm{Type}}_{i+1}$
are axioms. Adapting the presentation to a hierarchy of sorts is
straightforward, but in the interest of keeping to the heart of the matter we
give a presentation with two sorts.
## 3\. Inductive type families
We shall now extend the algebraic presentation with a notion of _inductive
type families_ to model (dependent) datatypes. In this section we will stay in
the predicative fragment of Coq.
Contrary to the inductive types of Coq, where inductive definitions must be
_named_ at _toplevel_ , like in:
Inductive Even : Type := ${\mid}$ eo : Even
${\mid}$ es : Odd ${\rightarrow}$ Even
with Odd : Type := ${\mid}$ os : Even ${\rightarrow}$ Odd.
the presentation given here is essentially anonymous, and inductive families
need not be defined at toplevel prior to use. Mutually inductive types such as
Even and Odd are not modelled directly, rather they are encoded using an
index:
Inductive EvenOdd : bool ${\rightarrow}$ Type := ${\mid}$ eo : EvenOdd true
${\mid}$ es : EvenOdd false ${\rightarrow}$ EvenOdd true
${\mid}$ os : EvenOdd true ${\rightarrow}$ EvenOdd false. Definition Even :=
EvenOdd true. Definition Odd := EvenOdd false.
This encoding works as long as all the mutual definitions are all in the same
sort. A variant for mutual definition involving Type and Prop is demonstrated
in Section 4.3. When the types being defined are in different predicative
sorts, however, we have to resort to another encoding which involves nested
datatypes [4, Section 8.6].
We will not explore the latter kind of mutual definition. However, nested
datatypes – where recursive calls occur as arguments of another type – such
as:
Inductive List (A:Type) : Type := ${\mid}$ nil : List A
${\mid}$ cons : A ${\rightarrow}$ List A ${\rightarrow}$ List A. Inductive
Tree : Type := ${\mid}$ node : List Tree ${\rightarrow}$ Tree.
are indeed modelled in the algebraic presentation.
### 3.1. Regular types
To be able to traverse terms of inductive type, the core pts constructions is
extended with a _recursive fixed point_ on functions:
$\Gamma\vdash\mathrm{fix}\leavevmode\nobreak\ f\leavevmode\nobreak\
{x}_{1}{:}{A}_{1}\leavevmode\nobreak\ \mbox{\ldots}\leavevmode\nobreak\
{x}_{n-1}{:}{A}_{n-1}\leavevmode\nobreak\ {x}_{n}{:}{A}_{n}\Rightarrow
u:\prod_{{x}_{1}{:}{A}_{1},\mbox{\ldots},{x}_{n-1}{:}{A}_{n-1},{x}_{n}{:}{A}_{n}}B\begin{array}[]{c}\eqmakebox[tag1]{$\Gamma\vdash\prod_{{x}_{1}{:}{A}_{1},\mbox{\ldots},{x}_{n-1}{:}{A}_{n-1},{x}_{n}{:}{A}_{n}}B
:
s$${\qquad}$${\mathrm{guarded}}\,{f}\,{{x}_{1}}\,{\mbox{\ldots}}\,{{x}_{n-1}}\,{{x}_{n}}\,{u}$}\\\
\eqmakebox[tag1]{$\Gamma,
f{:}\prod_{{x}_{1}{:}{A}_{1},\mbox{\ldots},{x}_{n-1}{:}{A}_{n-1},{x}_{n}{:}{A}_{n}}B,
{x}_{1}{:}{A}_{1}, \mbox{\ldots}, {x}_{n-1}{:}{A}_{n-1},
{x}_{n}{:}{A}_{n}\vdash u : B$}\\\ \end{array}$
Recursive fixed points are unfolded when fully applied
${\left(\mathrm{fix}\leavevmode\nobreak\ f\leavevmode\nobreak\
{x}_{1}{:}{A}_{1}\leavevmode\nobreak\ \mbox{\ldots}\leavevmode\nobreak\
{x}_{n}{:}{A}_{n}\Rightarrow
u\right)}\,{{v}_{1}}\,{\mbox{\ldots}}\,{{v}_{n}}\leadsto{u}{\left[\,{x}_{i}\,{\setminus}\,{v}_{i}\,\right]}$
To ensure strong normalisation, this reduction rule is limited, and a guard
condition is imposed on the recursive calls to $f$. It is not, however, the
object of this article to discuss these restriction or the guard condition.
Briefly, Coq currently relies on a single _structural_ argument in the block
${x}_{1}$, …, ${x}_{n}$: fixed points are not unfolded until their structural
argument starts with a constructor, and the guard condition ensures that each
recursive call is performed on a subterm of said structural argument, for a
relaxed notion of subterm. Other possibilities exist: Agda2 [11] uses any
number of arguments as structural, and tries to find a lexicographic ordering.
Yet another possibilities is to use sized types [1]. We shall simply assume
that an adequate guard condition is given.
We now extend the grammar of type constructors. The presentation of this
article is largely inspired by the synthetic definition of _strictly positive
families_ by Morris & al [9, 10], but is adapted to _intensional type theory_.
The presentation of [9, 10] is designed for generic programming inside a type
theory, they give codes for strictly positive families which are then decoded
into an actual type of the ambient theory. No elimination principle needs to
be given for the strictly positive families, as they are implicit in their
decoding. Here, we are defining the syntax of inductive type families,
including their elimination rules.
Sum type
$\Gamma\vdash A+B:\mathrm{Type}\lx@proof@logical@and\Gamma\vdash
A:\mathrm{Type}\Gamma\vdash B:\mathrm{Type}$
$\Gamma\vdash{\mathrm{inl}}\,{t}:A+B\lx@proof@logical@and\Gamma\vdash
A+B:s\Gamma\vdash t:A$
$\Gamma\vdash{\mathrm{inr}}\,{u}:A+B\lx@proof@logical@and\Gamma\vdash
A+B:s\Gamma\vdash u:B$
$\Gamma\vdash\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{inl}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\
{\mathrm{inr}}\,{z}\hfil\leavevmode\nobreak\ \Rightarrow\leavevmode\nobreak\
&w\\\ \end{array}\\\
\end{array}:{P}{\left[\,x\,{\setminus}\,u\,\right]}\lx@proof@logical@and\Gamma\vdash
u:A+B\Gamma,x{:}A+B\vdash P:s\Gamma,y{:}A\vdash
v:{P}{\left[\,x\,{\setminus}\,{\mathrm{inl}}\,{y}\,\right]}\Gamma,z{:}A\vdash
w:{P}{\left[\,x\,{\setminus}\,{\mathrm{inr}}\,{z}\,\right]}$
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
{\mathrm{inl}}\,{u}\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\
x\leavevmode\nobreak\ \mathrm{return}\leavevmode\nobreak\
P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{inl}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\
{\mathrm{inr}}\,{z}\hfil\leavevmode\nobreak\ \Rightarrow\leavevmode\nobreak\
&w\\\ \end{array}\\\
\end{array}\leadsto{v}{\left[\,y\,{\setminus}\,u\,\right]}$
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
{\mathrm{inr}}\,{u}\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\
x\leavevmode\nobreak\ \mathrm{return}\leavevmode\nobreak\
P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{inl}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\
{\mathrm{inr}}\,{z}\hfil\leavevmode\nobreak\ \Rightarrow\leavevmode\nobreak\
&w\\\ \end{array}\\\
\end{array}\leadsto{w}{\left[\,z\,{\setminus}\,u\,\right]}$
Unit type
$\Gamma\vdash 1:\mathrm{Type}$ $\Gamma\vdash():1$
$\Gamma\vdash\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}()\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}:{P}{\left[\,x\,{\setminus}\,u\,\right]}\lx@proof@logical@and\Gamma\vdash
u:1\Gamma,x{:}1\vdash P:s\Gamma\vdash
v:{P}{\left[\,x\,{\setminus}\,()\,\right]}$
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
()\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}()\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\ \end{array}\leadsto v$
Empty type
$\Gamma\vdash 0:\mathrm{Type}$ $\Gamma\vdash\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{return}\leavevmode\nobreak\
A\leavevmode\nobreak\ \mathrm{with}\leavevmode\nobreak\
\cdot:A\lx@proof@logical@and\Gamma\vdash u:0\Gamma\vdash A:\mathrm{Type}$
Figure 2. Regular type constructors
The regular type constructors, whose typing rules are given in Figure 2, are
the empty type $0$, the unit type $1$, and the sum of two types. The
elimination rules are given in the form of _dependent pattern-matching_ with a
syntax made to remind of that of Coq. We shall often omit the typing predicate
when it is clear from the context, especially when it does not depend on the
branch. With this material we can define a first example type, namely the
booleans:
$\left\\{\begin{aligned} \mathbb{B}&=1+1\\\
\mathrm{true}&={\mathrm{inl}}\,{()}\\\
\mathrm{false}&={\mathrm{inr}}\,{()}\end{aligned}\right.$
### 3.2. Inductive type families
Inductive families differ from regular inductive types in that they are
parametrised by _indices_ , that is they are functions
$F:A\rightarrow\mathrm{Type}$ for some $A$. An inductive family of the form
$\lambda{x}^{A}.\,R$, is said to be _uniformly parametrised_ by $A$. In
general, inductive families are not uniformly parametrised: the value of the
index is allowed to vary in recursive calls, and constructors may build values
of ${F}\,{x}$ for certain $x$ only. Remember, for instance, the EvenOdd
family:
Inductive EvenOdd : bool ${\rightarrow}$ Type := ${\mid}$ eo : EvenOdd true
${\mid}$ es : EvenOdd false ${\rightarrow}$ EvenOdd true
${\mid}$ os : EvenOdd true ${\rightarrow}$ EvenOdd false.
The inductive family constructors, presented in Figure 3, warrant individual
discussion. First, notice that as a simplifying hypothesis, inductive families
have exactly one index. This is, of course, not a limitation in expressive
power as multiple indices can be encoded as a dependent sum, and the unit type
allows us to write families without an index.
Inductive fixed point
$\Gamma\vdash\mu{X}^{A\rightarrow\mathrm{Type}}.\,F:A\rightarrow\mathrm{Type}\lx@proof@logical@and\Gamma\vdash
A:s\Gamma,X{:}A\rightarrow\mathrm{Type}\vdash
F:A\rightarrow\mathrm{Type}{{\mbox{{sp}}}_{X}}\,{F}$
$\mu{X}^{A\rightarrow
s}.\,F\equiv{F}{\left[\,X\,{\setminus}\,\mu{X}^{A\rightarrow s}.\,F\,\right]}$
Proper indices
$\Gamma\vdash{\sum}_{x{:}A}^{f}T:B\rightarrow\mathrm{Type}\lx@proof@logical@and\Gamma\vdash
A:\mathrm{Type}\Gamma\vdash B:\mathrm{Type}\Gamma,x{:}A\vdash
T:\mathrm{Type}\Gamma,x{:}A\vdash f:B$
$\Gamma\vdash{\left(u,v\right)}_{x{:}A.T}^{f}:{\left({\sum}_{x{:}A}^{f}T\right)}\,{\left({f}{\left[\,x\,{\setminus}\,u\,\right]}\right)}\lx@proof@logical@and\Gamma\vdash{\sum}_{x{:}A}^{f}T:B\rightarrow
s\Gamma\vdash u:A\Gamma\vdash v:{T}{\left[\,x\,{\setminus}\,u\,\right]}$
$\Gamma\vdash\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ z\leavevmode\nobreak\
\mathrm{in}\leavevmode\nobreak\ y\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\left(i,j\right)}_{x{:}A.T}^{f}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}:{P}{\left[\,y\,{\setminus}\,b\,,\,z\,{\setminus}\,u\,\right]}\begin{array}[]{c}\eqmakebox[tag2]{$\Gamma\vdash
u : {\left( {\sum}_{x{:}A}^{f}T\right) }\,{b}$${\qquad}$$\Gamma, y{:}B,
z{:}{\left( {\sum}_{x{:}A}^{f}T\right) }\,{y}\vdash P : s$}\\\
\eqmakebox[tag2]{$\Gamma, i{:}A, j{:}{T}\,{i}\vdash v : {P}{\left[
\,y\,{\setminus}\,{f}\,{i}\,,\,z\,{\setminus}\,{\left( i, j\right)
}_{x{:}A.T}^{f}\,\right] }$}\\\ \end{array}$
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
{\left(u,v\right)}_{x{:}A.T}^{f}\leavevmode\nobreak\
\mathrm{as}\leavevmode\nobreak\ z\leavevmode\nobreak\
\mathrm{in}\leavevmode\nobreak\ y\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\left(i,j\right)}_{x{:}A.T}^{f}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &w\\\ \end{array}\\\
\end{array}\leadsto{w}{\left[\,i\,{\setminus}\,u\,,\,j\,{\setminus}\,v\,\right]}$
Figure 3. Inductive type families
${{\mbox{{sp}}}_{X}}\,{y}$ ${{\mbox{{sp}}}_{X}}\,{0}$
${{\mbox{{sp}}}_{X}}\,{1}$
${{\mbox{{sp}}}_{X}}\,{\left(\prod_{x{:}A}B\right)}\lx@proof@logical@and
X\mbox{{ is fresh in }}A{{\mbox{{sp}}}_{X}}\,{B}$
${{\mbox{{sp}}}_{X}}\,{\left(A+B\right)}\lx@proof@logical@and{{\mbox{{sp}}}_{X}}\,{A}{{\mbox{{sp}}}_{X}}\,{B}$
${{\mbox{{sp}}}_{X}}\,{\left(\mu{Y}^{A}.\,F\right)}\lx@proof@logical@and
X\mbox{{ is fresh in }}A{{\mbox{{sp}}}_{X}}\,{F}$
${{\mbox{{sp}}}_{X}}\,{\left(\lambda{x}^{A}.\,T\right)}\lx@proof@logical@and
X\mbox{{ is fresh in }}A{{\mbox{{sp}}}_{X}}\,{T}$
${{\mbox{{sp}}}_{X}}\,{\left({U}\,{t}\right)}\lx@proof@logical@and{{\mbox{{sp}}}_{X}}\,{U}X\mbox{{
is fresh in }}t$
${{\mbox{{sp}}}_{X}}\,{A}\lx@proof@logical@and{{\mbox{{sp}}}_{X}}\,{B}B\equiv
A$
${{\mbox{{sp}}}_{X}}\,{\left({\sum}_{x{:}A}^{f}T\right)}\lx@proof@logical@and
X\mbox{{ is fresh in }}f{{\mbox{{sp}}}_{X}}\,{A}{{\mbox{{sp}}}_{X}}\,{T}$
Figure 4. Strict positivity condition
The construction $\mu{X}^{A\rightarrow\mathrm{Type}}.\,F$ constructs the
_inductive fixed point_ of $F$. It acts on type families, because indices vary
through recursive calls to $X$. To be able to form an inductive fixed point,
occurrences of $X$ must be strictly positive in $F$, rules for strict
positivity are given in Figure 4. The rules of Figure 4 are a simple set which
suits the needs of this article, however in practice, we may want to consider
strict positivity rules involving elimination rules and a finer treatment of
application. Strict positivity ensures that no non-terminating term can be
written without recursive fixed points, so that the guard condition suffices
to enforce termination. Paradoxes which can be derived from non-positive or
non-strictly positive inductive fixed points can be found in [13, Chapter 4,
Section 4.2][8, Chapter 3][4, Chapter 8]. To avoid clutter, we give a
presentation where inductive fixed points can be freely rolled and unrolled
thanks to the conversion. An alternative can be to give an explicit term
constructor for fixed points, see Section 3.4.
We will also use an inductive fixed point on nullary families, defined as:
$\displaystyle\mu{X}^{\mathrm{Type}}.\,F$
$\displaystyle={\left(\mu{Y}^{1\rightarrow\mathrm{Type}}.\,{F}{\left[\,X\,{\setminus}\,{Y}\,{()}\,\right]}\right)}\,{()}$
from which we have that $\mu{X}^{\mathrm{Type}}.\,F$ can be freely rolled from
or unrolled to
${F}{\left[\,X\,{\setminus}\,\mu{X}^{\mathrm{Type}}.\,F\,\right]}$.
With inductive fixed points, we can, for instance, define the accessibility
predicate. In Coq:
Inductive Acc (A:Type) (R:A${\rightarrow}$A${\rightarrow}$Type) (x:A) : Type
:= ${\mid}$ acc_intro : (forall y:A, R y x ${\rightarrow}$ Acc A R y)
${\rightarrow}$ Acc A R x.
This type represents the generic form of termination proofs: any terminating
recursive fixed point can be made structural over a proof of accessibility. In
the algebraic presentation, it is defined as:
$\left\\{\begin{aligned}
\mathrm{Acc}&=\lambda{A}^{\mathrm{Type}}\,{R}^{A\rightarrow
A\rightarrow\mathrm{Type}}.\,\mu{\mathrm{Acc}}^{A\rightarrow\mathrm{Type}}.\,\lambda{x}^{A}.\,\prod_{y{:}A}{R}\,{y}\,{x}\rightarrow{\mathrm{Acc}}\,{y}\\\
{\mathrm{acc}}_{\mathrm{intro}}&=\lambda{A}^{\mathrm{Type}}\,{R}^{A\rightarrow
A\rightarrow\mathrm{Type}}\,{x}^{A}\,{f}^{\prod_{y{:}A}{R}\,{y}\,{x}\rightarrow{\mathrm{Acc}}\,{A}\,{R}\,{y}}.\,f\end{aligned}\right.$
Because inductive fixed points are treated transparently, the constructor is
rather trivial. However, notice how, in the definition of $\mathrm{Acc}$, the
parameter $x$ is treated differently from $A$ and $R$. The reason is that $A$
and $R$ are _uniform parameters_ , in that they do not vary through recursive
calls, whereas $x$ does: it is a _non-uniform parameter_. The parameter $x$
is, hence, represented as an index. However, such an index is not sufficient
to encode types like EvenOdd.
Representing proper indices requires a new type construction, which we write
${\sum}_{x{:}A}^{f}T$. This construction comes from [9, 10], where it is
inspired by a categorical point of view: in a sufficiently extensional
setting, ${\sum}_{x{:}A}^{f}T$ is the right adjoint to a pullback functor. The
similarity with the usual notation of dependent sum is not fortuitous, indeed
we can define dependent sum as a special case of proper indexing:
$\left\\{\begin{aligned}
\sum_{x{:}A}B&={\left({\sum}_{x{:}A}^{()}B\right)}\,{()}\\\
{\left(u,v\right)}_{{\mbox{\scriptsize{$x${:}$A$.$B$}}}}&={\left(u,v\right)}_{x{:}A.B}^{()}\end{aligned}\right.$
We also write $A\times B$ and $\left(u,v\right)$ for $\sum_{x{:}A}B$ and
${\left(u,v\right)}_{{\mbox{\scriptsize{$x${:}$A$.$B$}}}}$, respectively, when
$x$ is not free in $B$.
In the case of dependent sums, the index is trivial. When it is not, however,
the pattern matching return clause $P$ is allowed to depend on the value of
the index. This is the purpose of Coq’s in-pattern. With the algebraic
presentation, the in-pattern has the pleasant property of being confined to
the proper indexing construction, hopefully making its meaning more explicit.
The syntax differs a little from that of Coq, however: Coq renders the in
clause as a pattern with the type name at the head:
match n as n’ in EvenOdd b return P n’ b with
… end.
In the algebraic presentation, types not having a name, the in clause simply
consists of a name for the index.
The prototype of proper indexing is the identity type, which we name
$\mathrm{Eq}$. In Coq:
Inductive Eq (A:Type) (x:A) : A ${\rightarrow}$ Type := ${\mid}$ eq_refl : Eq
A x x
in the algebraic presentation:
$\left\\{\begin{aligned}
\mathrm{Eq}&=\lambda{A}^{\mathrm{Type}}\,{x}^{A}.\,{\sum}_{\\_{:}1}^{x}1\\\
{\mathrm{eq}}_{\mathrm{refl}}&=\lambda{A}^{\mathrm{Type}}\,{x}^{A}.\,{\left((),()\right)}_{\\_{:}1.1}^{x}\end{aligned}\right.$
In fact, dependent sums and identity types are sufficient to define proper
indexing. Indeed ${\sum}_{x{:}A}^{f}T$ can be redefined as:
$\left\\{\begin{aligned}
{\sum}_{x{:}A}^{f}T&=\lambda{y}^{B}.\,\sum_{x{:}A}\left({\mathrm{Eq}}\,{B}\,{y}\,{f}\right)\times
T\\\
{\left(u,v\right)}_{x{:}A.T}^{f}&={\left(u,\left({{\mathrm{eq}}_{\mathrm{refl}}}\,{B}\,{{f}{\left[\,x\,{\setminus}\,u\,\right]}},v\right)\right)}_{{\mbox{\scriptsize{$x${:}$A$.$\left({\mathrm{Eq}}\,{B}\,{\left({f}{\left[\,x\,{\setminus}\,u\,\right]}\right)}\,{f}\right)$${\times}$$T$}}}}\end{aligned}\right.$
It is closer to the spirit of Coq, but in no way essential, to take a proper
indexing construction rather than equality as primitive. In Morris & al [9,
10], the dependent sum and equality of the ambient type theory is used to
define ${\sum}_{x{:}A}^{f}T$ which is then taken as primitive.
An other choice lies in the use of $A+B$ as primitive. It is the only type
construction which allows to define a type with distinct elements. However, a
common alternative is to take $\mathbb{B}$ as primitive, in which case we can
define $A+B$ as:
$\left\\{\begin{aligned}
A+B&=\sum_{b{:}\mathbb{B}}\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
b\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}\mathrm{true}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &A\\\ \mathrm{false}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &B\\\ \end{array}\\\ \end{array}\\\
\mathrm{inl}&=\lambda{x}^{A}.\,{\left(\mathrm{true},x\right)}_{{\mbox{\scriptsize{$b${:}$\mathbb{B}$.$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
b\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}\mathrm{true}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &A\\\ \mathrm{false}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &B\\\ \end{array}\\\ \end{array}$}}}}\\\
\mathrm{inr}&=\lambda{y}^{B}.\,{\left(\mathrm{false},y\right)}_{{\mbox{\scriptsize{$b${:}$\mathbb{B}$.$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
b\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}\mathrm{true}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &A\\\ \mathrm{false}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &B\\\ \end{array}\\\
\end{array}$}}}}\end{aligned}\right.$
### 3.3. Examples
The previous section concludes the description of the predicative fragment of
the algebraic presentation. We can now give definitions of the remaining
inductive families we have encountered in terms of the algebraic presentation.
Starting with EvenOdd:
Inductive EvenOdd : bool ${\rightarrow}$ Type := ${\mid}$ eo : EvenOdd true
${\mid}$ es : EvenOdd false ${\rightarrow}$ EvenOdd true
${\mid}$ os : EvenOdd true ${\rightarrow}$ EvenOdd false.
translates, in the algebraic presentation, to:
$\left\\{\begin{aligned}
\mathrm{EvenOdd}&=\mu{\mathrm{EvenOdd}}^{\mathbb{B}\rightarrow\mathrm{Type}}.\,\left({\sum}_{\\_{:}1}^{\mathrm{true}}1+{\mathrm{EvenOdd}}\,{\mathrm{false}}\right)+\left({\sum}_{\\_{:}1}^{\mathrm{false}}{\mathrm{EvenOdd}}\,{\mathrm{true}}\right)\\\
\mathrm{eo}&={\mathrm{inl}}\,{{\left((),{\mathrm{inl}}\,{()}\right)}_{\\_{:}1.1}^{\mathrm{true}}}\\\
\mathrm{es}&=\lambda{x}^{{\mathrm{EvenOdd}}\,{\mathrm{false}}}.\,{\mathrm{inl}}\,{{\left((),{\mathrm{inr}}\,{x}\right)}_{\\_{:}1.1}^{\mathrm{true}}}\\\
\mathrm{os}&=\lambda{x}^{{\mathrm{EvenOdd}}\,{\mathrm{true}}}.\,{\mathrm{inr}}\,{{\left((),x\right)}_{\\_{:}1.1}^{\mathrm{false}}}\end{aligned}\right.$
Here is the definition of List in Coq:
Inductive List (A:Type) : Type := ${\mid}$ nil : List A
${\mid}$ cons : A ${\rightarrow}$ List A ${\rightarrow}$ List A.
and in the algebraic presentation:
$\left\\{\begin{aligned}
\mathrm{List}&=\lambda{A}^{\mathrm{Type}}.\,\mu{\mathrm{List}}^{\mathrm{Type}}.\,1+A\times\mathrm{List}\\\
\mathrm{nil}&=\lambda{A}^{\mathrm{Type}}.\,{\mathrm{inl}}\,{()}\\\
\mathrm{cons}&=\lambda{A}^{\mathrm{Type}}\,{x}^{A}\,{l}^{{\mathrm{List}}\,{A}}.\,{\mathrm{inr}}\,{\left(x,l\right)}\end{aligned}\right.$
Finally the definition of Tree:
Inductive Tree : Type := ${\mid}$ node : List Tree ${\rightarrow}$ Tree.
translates to:
$\left\\{\begin{aligned}
\mathrm{Tree}&=\mu{\mathrm{Tree}}^{\mathrm{Type}}.\,{\mathrm{List}}\,{\mathrm{Tree}}\\\
\mathrm{node}&=\lambda{f}^{{\mathrm{List}}\,{\mathrm{Tree}}}.\,f\end{aligned}\right.$
Note that in the definition of $\mathrm{Tree}$, we must ${\beta}$-reduce
${\mathrm{List}}\,{\mathrm{Tree}}$ for the recursive definition to be strictly
positive.
A more complex example is given by the type of lists indexed by their length,
often called _vectors_ :
Inductive Nat : Type := ${\mid}$ o : Nat
${\mid}$ s : Nat ${\rightarrow}$ Nat. Inductive Vector (A:Type) : Nat
${\rightarrow}$ Type := ${\mid}$ vnil : Vector A o
${\mid}$ vcons : forall n, A ${\rightarrow}$ Vector A n ${\rightarrow}$ Vector
A (s n).
It is encoded in the algebraic presentation as:
$\left\\{\begin{aligned}
\mathrm{Nat}&=\mu{\mathrm{Nat}}^{\mathrm{Type}}.\,1+\mathrm{Nat}\\\
o&={\mathrm{inl}}\,{()}\\\
s&=\lambda{n}^{\mathrm{Nat}}.\,{\mathrm{inr}}\,{n}\\\
\mathrm{Vector}&=\lambda{A}^{\mathrm{Type}}.\,\mu{V}^{\mathrm{Nat}\rightarrow\mathrm{Type}}.\,\lambda{n}^{\mathrm{Nat}}.\,{\left({\sum}_{\\_{:}1}^{o}1\right)}\,{n}+{\left({\sum}_{n^{\prime}{:}\mathrm{Nat}}^{{s}\,{n^{\prime}}}A\times{V}\,{n^{\prime}}\right)}\,{n}\\\
\mathrm{vnil}&=\lambda{A}^{\mathrm{Type}}.\,{\mathrm{inl}}\,{{\left((),()\right)}_{\\_{:}1.1}^{o}}\\\
\mathrm{vcons}&=\lambda{A}^{\mathrm{Type}}\,{n}^{\mathrm{Nat}}\,{a}^{A}\,{v}^{{\mathrm{Vector}}\,{n}\,{a}}.\,{\mathrm{inr}}\,{{\left(n,\left(a,v\right)\right)}_{n^{\prime}{:}\mathrm{Nat}.A\times{V}\,{n^{\prime}}}^{{s}\,{n^{\prime}}}}\end{aligned}\right.$
Contrary to proper indices, the types of non-uniform parameters are allowed to
be in $\Box$, this allows the definition of types such as the binary lists
[12]:
Inductive BList (A:Type) : Type := ${\mid}$ one : A ${\rightarrow}$ BList A
${\mid}$ twice : BList (A$*$A) ${\rightarrow}$ BList A
${\mid}$ stwice : A ${\rightarrow}$ BList (A$*$A) ${\rightarrow}$ BList A
which are rendered in the algebraic presentation as:
$\left\\{\begin{aligned}
\mathrm{BList}&=\mu{\mathrm{BList}}^{\mathrm{Type}\rightarrow\mathrm{Type}}.\,\lambda{A}^{\mathrm{Type}}.\,A+\left({\mathrm{BList}}\,{\left(A\times
A\right)}+A\times\left({\mathrm{BList}}\,{\left(A\times
A\right)}\right)\right)\\\
\mathrm{one}&=\lambda{A}^{\mathrm{Type}}\,{x}^{A}.\,{\mathrm{inl}}\,{x}\\\
\mathrm{twice}&=\lambda{A}^{\mathrm{Type}}\,{l}^{{\mathrm{BList}}\,{\left(A\times
A\right)}}.\,{\mathrm{inr}}\,{\left({\mathrm{inl}}\,{l}\right)}\\\
\mathrm{stwice}&=\lambda{A}^{\mathrm{Type}}\,{x}^{A}\,{l}^{{\mathrm{BList}}\,{\left(A\times
A\right)}}.\,{\mathrm{inr}}\,{\left({\mathrm{inr}}\,{\left(a,l\right)}\right)}\end{aligned}\right.$
In Coq, where there is a hierarchy of universe, types of proper indices can be
in any sort. However, a proper index whose type is in ${\mathrm{Type}}_{i}$
constrains the final type to be in ${\mathrm{Type}}_{i+1}$ or higher. Uniform
parameters, of any type, do not constrain the type they parametrise.
Inductive types are consumed by recursive fixed points. Using the implicit
unfolding of inductive fixed points, we can pattern-match over the top
constructor directly. The Coq function
Fixpoint add (x y:Nat) : Nat := match y with
${\mid}$ o ${\Rightarrow}$ x
${\mid}$ s y’ ${\Rightarrow}$ s (add x y’) end
is rendered in the algebraic presentation as
$\displaystyle\mathrm{add}$ $\displaystyle=\mathrm{fix}\leavevmode\nobreak\
\mathrm{add}\leavevmode\nobreak\ x{:}\mathrm{Nat}\leavevmode\nobreak\
y{:}\mathrm{Nat}\Rightarrow\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
y\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ \\_\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ \mathrm{Nat}\leavevmode\nobreak\
\mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{inl}}\,{\\_}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &x\\\
{\mathrm{inr}}\,{y^{\prime}}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\
&{s}\,{\left({\mathrm{add}}\,{x}\,{y^{\prime}}\right)}\\\ \end{array}\\\
\end{array}$
### 3.4. Co-induction
In addition to inductive fixed points, Coq also has support for co-inductive
fixed points. Co-inductive fixed points are required to be strictly positive,
like inductive fixed points. We choose in this section, a presentation of co-
inductive data where fixed points are explicitly introduced with a
constructor. Below we will use this explicit presentation to give a variation
on Coq’s co-inductive fixed points.
$\parbox{433.62pt}{ \begin{center} \makebox[433.62pt]{\parbox{433.62pt}{
\begin{center}
$\Gamma\vdash\nu{X}^{A\rightarrow\mathrm{Type}}.\,F:A\rightarrow\mathrm{Type}\lx@proof@logical@and\Gamma\vdash
A:s\Gamma,X{:}A\rightarrow\mathrm{Type}\vdash
F:A\rightarrow\mathrm{Type}{{\mbox{{sp}}}_{X}}\,{F}$ \end{center} }}\\\
\makebox[433.62pt]{\parbox{433.62pt}{ \begin{center}
$\Gamma\vdash{\mathrm{forced}}\,{u}:{\left(\nu{X}^{A\rightarrow
s}.\,F\right)}\,{i}\lx@proof@logical@and\Gamma\vdash\nu{X}^{A\rightarrow
s}.\,F:A\rightarrow s\Gamma\vdash i:A\Gamma\vdash
u:{{F}{\left[\,X\,{\setminus}\,\nu{X}^{A\rightarrow s}.\,F\,\right]}}\,{i}$
\end{center} }}\\\ \makebox[433.62pt]{\parbox{433.62pt}{ \begin{center}
$\Gamma\vdash\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{forced}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}:{P}{\left[\,x\,{\setminus}\,u\,\right]}\begin{array}[]{c}\eqmakebox[tag3]{$\Gamma\vdash\nu{X}^{A\rightarrow
s}.\,F : A\rightarrow s$${\qquad}$$\Gamma\vdash i : A$${\qquad}$$\Gamma\vdash
u : {\left( \nu{X}^{A\rightarrow s}.\,F\right) }\,{i}$}\\\
\eqmakebox[tag3]{$\Gamma, x{:}{\left( \nu{X}^{A\rightarrow s}.\,F\right)
}\,{i}\vdash P : s'$${\qquad}$$\Gamma, y{:}{{F}{\left[
\,X\,{\setminus}\,\nu{X}^{A\rightarrow s}.\,F\,\right] }}\,{i}\vdash v :
{P}{\left[ \,x\,{\setminus}\,{\mathrm{forced}}\,{y}\,\right] }$}\\\
\end{array}$ \end{center} }}\\\ \makebox[433.62pt]{\parbox{433.62pt}{
\begin{center} $\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
{\mathrm{forced}}\,{u}\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\
x\leavevmode\nobreak\ \mathrm{return}\leavevmode\nobreak\
P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{forced}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}\leadsto{v}{\left[\,y\,{\setminus}\,u\,\right]}$ \end{center} }}\\\
\end{center} }$
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
{\mathrm{forced}}\,{u}\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\
x\leavevmode\nobreak\ \mathrm{return}\leavevmode\nobreak\
P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{forced}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}\leadsto{v}{\left[\,y\,{\setminus}\,u\,\right]}$
Just like inductive data is _destructed_ by a recursive fixed point operation,
co-inductive data is _constructed_ by a co-recursive fixed point operation,
allowing co-inductive data to be infinite. The guard condition on co-recursive
fixed points ensures that a finite number of unfolding will eventually produce
a forced value.
$\Gamma\vdash\mathrm{cofix}\leavevmode\nobreak\ f\leavevmode\nobreak\
{x}_{1}{:}{A}_{1}\leavevmode\nobreak\ \mbox{\ldots}\leavevmode\nobreak\
{x}_{n}{:}{A}_{n}\Rightarrow
u:\prod_{{x}_{1}{:}{A}_{1},\mbox{\ldots},{x}_{n}{:}{A}_{n}}{\left(\nu{X}^{A\rightarrow
s}.\,F\right)}\,{i}\begin{array}[]{c}\eqmakebox[tag4]{$\Gamma\vdash\prod_{{x}_{1}{:}{A}_{1},\mbox{\ldots},{x}_{n}{:}{A}_{n}}{\left(
\nu{X}^{A\rightarrow s}.\,F\right) }\,{i} :
s$${\qquad}$${\mathrm{coguarded}}\,{f}\,{{x}_{1}}\,{\mbox{\ldots}}\,{{x}_{n}}\,{u}$}\\\
\eqmakebox[tag4]{$\Gamma,
f{:}\prod_{{x}_{1}{:}{A}_{1},\mbox{\ldots},{x}_{n}{:}{A}_{n}}{\left(
\nu{X}^{A\rightarrow s}.\,F\right) }\,{i}, {x}_{1}{:}{A}_{1}, \mbox{\ldots},
{x}_{n}{:}{A}_{n}\vdash u : {\left( \nu{X}^{A\rightarrow s}.\,F\right)
}\,{i}$}\\\ \end{array}$
Co-recursive fixed-point are meant to represent infinite data: they cannot be
unfolded eagerly, lest they would fail to terminate. They are unfolded only
when they appear at the head of a pattern-matching expression:
$\begin{array}[]{ll}&\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
{\left(\mathrm{cofix}\leavevmode\nobreak\ f\leavevmode\nobreak\
{x}_{1}{:}{A}_{1}\leavevmode\nobreak\ \mbox{\ldots}\leavevmode\nobreak\
{x}_{n}{:}{A}_{n}\Rightarrow
u\right)}\,{{v}_{1}}\,{\mbox{\ldots}}\,{{v}_{n}}\leavevmode\nobreak\
\mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{forced}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\ \end{array}\\\
\leadsto&\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
{u}{\left[\,f\,{\setminus}\,\left(\mathrm{cofix}\leavevmode\nobreak\
f\leavevmode\nobreak\ {x}_{1}{:}{A}_{1}\leavevmode\nobreak\
\mbox{\ldots}\leavevmode\nobreak\ {x}_{n}{:}{A}_{n}\Rightarrow
u\right)\,,\,{x}_{i}\,{\setminus}\,{v}_{i}\,\right]}\leavevmode\nobreak\
\mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{forced}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\ \end{array}\\\
\end{array}$
The dependent elimination rule for co-inductive fixed points asserts, in
essence, that every co-inductive data is of the form ${\mathrm{forced}}\,{u}$.
Even though it would be fine for inductive fixed points – this is why we could
leave the unrolling to the conversion – this does not reflect well the
computational aspects of co-inductive data: suspended co-recursive fixed
points are values, and won’t be evaluated until the context demands it. The
fact that the elimination for co-inductive data claims that all values are
forced gives rise to undesirable behaviour.
Take for instance the following simple co-inductive type, and data:
$\left\\{\begin{aligned} T&=\nu X.\,X\\\ i&=\mathrm{cofix}\leavevmode\nobreak\
i\Rightarrow{\mathrm{forced}}\,{i}\end{aligned}\right.$
So that $i$ is effectively an infinite sequence of $\mathrm{forced}$. Using
the elimination principle above, it is possible to give a _closed_ proof that
${\mathrm{Eq}}\,{T}\,{i}\,{\left({\mathrm{forced}}\,{i}\right)}$:
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\ i\leavevmode\nobreak\
\mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\
{\mathrm{Eq}}\,{T}\,{x}\,{\left({\mathrm{forced}}\,{\left(\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
x\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{forced}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &y\\\ \end{array}\\\
\end{array}\right)}\right)}\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{forced}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &{{\mathrm{eq}}_{\mathrm{refl}}}\,{T}\,{y}\\\
\end{array}\\\ \end{array}$
However, $i$ and $\left({\mathrm{forced}}\,{i}\right)$ are not convertible,
yet, as every closed proof of equality does, this proof reduces to
${\mathrm{eq}}_{\mathrm{refl}}$, hence should relate convertible terms. The
dependent elimination rule of co-inductive fixed points compromises the type
safety of the logic.
Coq uses the above dependent elimination rule for co-inductive fixed points.
It was a deliberate decision made for practical purposes. Nonetheless, one may
want to weaken it to avoid the incompatibility between equality and
conversion. To do so, it suffices to erase the dependency of the return
predicate over the matched term:
$\Gamma\vdash\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{return}\leavevmode\nobreak\
P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{forced}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}:P\lx@proof@logical@and\Gamma\vdash\nu{X}^{A\rightarrow
s}.\,F:A\rightarrow s\Gamma\vdash i:A\Gamma\vdash
u:{\left(\nu{X}^{A\rightarrow s}.\,F\right)}\,{i}\Gamma\vdash
P:s^{\prime}\Gamma\vdash v:P$
## 4\. Prop
With all the common baggage for predicative sorts set in place, we can add
impredicative sorts to the algebraic presentation. The main such sort in Coq
is the sort Prop of propositions. The design of Prop is guided by _proof
irrelevance_ : even if it is not actually provable in Coq, different proofs of
a proposition are thought of as being equal. This property is useful for
program extraction: only the _computationally relevant_ parts of a program
need to be executed to get the final result. In other words: propositions are
considered as _static_ data. It is why, with disjunction and existential
defined as:
Inductive Or (A B:Prop) : Prop := ${\mid}$ or_introl : A ${\rightarrow}$ A
${\lor}$ B
${\mid}$ or_intror : B ${\rightarrow}$ A ${\lor}$ B. Inductive Ex (A:Type)
(P:A${\rightarrow}$Prop) : Prop := ${\mid}$ ex_intro : forall x:A, P x
${\rightarrow}$ Ex A P.
the following terms are refused by type-checking:
match x with
${\mid}$ or_introl $\\_$ ${\Rightarrow}$ true
${\mid}$ or_intror $\\_$ ${\Rightarrow}$ false
end.
and
match x with
${\mid}$ ex_intro x $\\_$ ${\Rightarrow}$ x
end.
On the other hand, it is not the case of every inductive type defined in Prop,
that they cannot be eliminated into Type. Conjunction and falsity are two
counter-examples:
Inductive False : Prop := . Inductive And (A B:Prop) : Prop := ${\mid}$ conj :
A ${\rightarrow}$ B ${\rightarrow}$ A ${\land}$ B.
Coq allows elimination over these two propositions into Type, and both
following terms are well-typed:
match x return Bool with end.
and
match x with
${\mid}$ conj $\\_$ $\\_$ ${\Rightarrow}$ true
end.
The object of this section is to make syntactically explicit what happens when
an inductive type of Coq is declared to be of sort Prop. The description
elaborated in this section has strong similarities with the system of bracket-
types proposed by Awodey & Bauer [2]. They describe the propositions as the
subset of types with at most one element, and introduce a left adjoint,
written as brackets, to the inclusion of propositions into types. We will
reuse their notation, even though, in our intensional setting, T:Prop does not
enforce that T has a most one element, and the bracketing operation does not
properly form an adjunction with the inclusion from Prop to Type.
### 4.1. Impredicativity
Let us start by introducing the new sort $\mathrm{Prop}$ in the algebraic
presentation:
$\Gamma\vdash\mathrm{Prop}:\mathrm{Type}$
As in [2], propositions form a subset of types. Coq has a subtyping rule (also
known as _cumulativity_) to make the inclusion transparent. We will, however,
render it with a syntactic construct:
$\Gamma\vdash\left\\{A\right\\}:\mathrm{Type}\Gamma\vdash A:\mathrm{Prop}$
$\Gamma\vdash{\mathrm{prf}}\,{u}:\left\\{A\right\\}\Gamma\vdash u:A$
$\Gamma\vdash\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{prf}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}:{P}{\left[\,x\,{\setminus}\,u\,\right]}\lx@proof@logical@and\Gamma\vdash
u:\left\\{A\right\\}\Gamma,x{:}\left\\{A\right\\}\vdash P:s\Gamma,y{:}A\vdash
v:{P}{\left[\,x\,{\setminus}\,{\mathrm{prf}}\,{y}\,\right]}$
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
{\mathrm{prf}}\,{u}\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\
x\leavevmode\nobreak\ \mathrm{return}\leavevmode\nobreak\
P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{prf}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}\leadsto{v}{\left[\,y\,{\setminus}\,u\,\right]}$
This definition simply makes $\left\\{A\right\\}$ a synonym of $A$, except of
sort $\mathrm{Type}$. It is strictly positive in $A$:
${{\mbox{{sp}}}_{X}}\,{\left\\{A\right\\}}{{\mbox{{sp}}}_{X}}\,{A}$
The fact that $\mathrm{Prop}$ is impredicative – _i.e._ supports the following
product formation rules:
$\Gamma\vdash\prod_{x{:}A}B:\mathrm{Prop}\lx@proof@logical@and\Gamma\vdash
A:s\Gamma,x{:}A\vdash B:\mathrm{Prop}$
is easily understood in terms of proof irrelevance. Indeed, if for all $x$,
$B$ has at most one element, so has the product over $x$. Even though it uses
functional extensionality, which is not provable.
### 4.2. Singleton rules
The types which (ideally) preserve the proof irrelevance property are
sometimes called singleton types in the setting of Coq. In our algebraic
presentation, they correspond to inductive type family constructors with extra
formation rules to make them preserve propositions. The rules are shown in
Figure 5.
$\Gamma\vdash 0:\mathrm{Prop}$ $\Gamma\vdash 1:\mathrm{Prop}$
$\Gamma\vdash{\sum}_{x{:}A}^{f}T:B\rightarrow\mathrm{Prop}\lx@proof@logical@and\Gamma\vdash
A:\mathrm{Prop}\Gamma\vdash B:\mathrm{Prop}\Gamma,x{:}A\vdash
T:\mathrm{Prop}\Gamma,x{:}A\vdash f:B$
Figure 5. Singleton rules
This elucidates why Coq allows elimination over False and And into arbitrary
type: False is implemented as $0$ and And A B and $A\times B$. The elimination
rules being unchanged, pattern-matching over proofs of False and And are
unrestricted. Because restricted pattern-matching is often seen as the
default, singleton types are said to enjoy _singleton elimination_.
Remark that, proofs of propositions being uninformative, there is essentially
nothing to be gained from depending on, or being indexed over a proposition.
In consequence, the type formation rule for proper indexing in Figure 5 is
only useful, in practice, for the subcase of cartesian product.
Coq actually implements two other singleton rules. The first one is for
inductive fixed points. In our algebraic presentation:
$\Gamma\vdash\mu{X}^{A\rightarrow\mathrm{Prop}}.\,F:A\rightarrow\mathrm{Prop}\lx@proof@logical@and\Gamma\vdash
A:s\Gamma,X{:}A\rightarrow\mathrm{Prop}\vdash
F:A\rightarrow\mathrm{Prop}{{\mbox{{sp}}}_{X}}\,{F}$
It allows to type the accessibility predicate $\mathrm{Acc}$ in
$\mathrm{Prop}$. This rule is sound in that fixed points indeed preserve proof
irrelevance in presence of functional extensionality. It is also very useful
for extraction: structural recursion over $\mathrm{Acc}$ allows the definition
of functions whose termination cannot be proved automatically by the guard
condition. However, the proof is no longer needed to ensure termination in the
target languages of extraction. In this sense, at least, it is static data.
The last singleton rule allows properly indexed families in $\mathrm{Prop}$
(not how it is stronger than the rule dependent sum of Figure 5):
$\Gamma\vdash{\sum}_{x{:}A}^{f}T:B\rightarrow\mathrm{Prop}\lx@proof@logical@and\Gamma\vdash
A:\mathrm{Prop}\Gamma\vdash B:\mathrm{Type}\Gamma,x{:}A\vdash
T:\mathrm{Prop}\Gamma,x{:}A\vdash f:B$
It turns the identity type $\mathrm{Eq}$ into a proposition. It is known to be
sound to accept that $\mathrm{Eq}$ is proof irrelevant [6]. It is also useful
for extraction, as equal types, in a closed environment, are extracted to the
same type. Hence a program may safely eliminate over $\mathrm{Eq}$ knowing
that it will not affect the performances of the extracted code. In Coq, the
index $B$ in the rule above can be of any sort ${\mathrm{Type}}_{i}$, however,
this wisdom has been challenged in recent years with the formulation of the
_univalence principle_ [15], of which a simple consequence is that
$\mathrm{Eq}$ is _not_ proof irrelevant at every type. Indeed, some extracted
Coq programs written assuming the univalence principle crash.
To correct for the univalence principle, the singleton rule for proper indices
can be simply dropped; but it can also be restricted to the lowest sort:
$B:{\mathrm{Type}}_{0}$. More precisely the conjunction of the univalence
principle and the proof irrelevance principle is consistent as long as the
singleton rule of proper indices is restricted to sorts $s$ such that there is
no sort $s^{\prime}$ other than $\mathrm{Prop}$ such that $s^{\prime}:s$.
Because, if such a sort $s^{\prime}$ exists, $\mathbb{B}:s^{\prime}$ and by
univalence, ${\mathrm{Eq}}\,{\mathbb{B}}$ has two distinct elements
contradicting proof irrelevance.
For types which do not enjoy singleton elimination, turning them into
propositions means restricting their elimination. We achieve this effect by
adding a single type construction coercing from $\mathrm{Type}$ to
$\mathrm{Prop}$:
$\Gamma\vdash\left[A\right]:\mathrm{Prop}\Gamma\vdash A:\mathrm{Type}$
$\Gamma\vdash\left\langle u\right\rangle:\left[A\right]\Gamma\vdash u:A$
$\Gamma\vdash\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}\left\langle
y\right\rangle\hfil\leavevmode\nobreak\ \Rightarrow\leavevmode\nobreak\ &v\\\
\end{array}\\\
\end{array}:{P}{\left[\,x\,{\setminus}\,u\,\right]}\lx@proof@logical@and\Gamma\vdash
u:\left[A\right]\mbox{${\quad}$$\Gamma,x{:}\left[A\right]\vdash
P:\mathrm{Prop}$${\quad}$}\Gamma,y{:}A\vdash
v:{P}{\left[\,x\,{\setminus}\,\left\langle y\right\rangle\,\right]}$
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\ \left\langle
u\right\rangle\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\
x\leavevmode\nobreak\ \mathrm{return}\leavevmode\nobreak\
P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}\left\langle
y\right\rangle\hfil\leavevmode\nobreak\ \Rightarrow\leavevmode\nobreak\ &v\\\
\end{array}\\\ \end{array}\leadsto{v}{\left[\,y\,{\setminus}\,u\,\right]}$
The important rule is the elimination rule, where the return clause is limited
to be of sort $\mathrm{Prop}$, whereas every other type construction can be
eliminated to any sort. Apart from this restriction $\left[A\right]$ is a
synonym of $A$, except in $\mathrm{Prop}$. In [2], the type theory is
extensional, in that the identity type and the conversion relation coincide.
The elimination rules for bracket is much finer and reflects precisely the
fact that propositions are proof-irrelevant. In an intensional type theory,
restricting with respect to sorts approximates this behaviour: even if we
constrained propositions to be proof-irrelevant, not every proof irrelevant
type will have type $\mathrm{Prop}$. The bracketing construction is also
strictly positive:
${{\mbox{{sp}}}_{X}}\,{\left[A\right]}{{\mbox{{sp}}}_{X}}\,{A}$
It is actually possible, using only the impredicative dependent product to
define a bracketing operation: $\prod_{P{:}\mathrm{Prop}}\left(A\rightarrow
P\right)\rightarrow P$. Like $\left[A\right]$ it behaves as $A$ except it can
only be used to form a proposition. However, the impredicative encoding is
positive but not strictly, which motivates the introduction of the extra
construction.
### 4.3. Examples
The logical connectives can be defined as follows:
$\left\\{\begin{aligned} \mathrm{False}&=0\\\
\mathrm{And}&=\lambda{A}^{\mathrm{Prop}}\,{B}^{\mathrm{Prop}}.\,A\times B\\\
\mathrm{pair}&=\lambda{A}^{\mathrm{Prop}}\,{B}^{\mathrm{Prop}}\,{x}^{A}\,{y}^{B}.\,\left(x,y\right)\\\
\mathrm{Or}&=\lambda{A}^{\mathrm{Prop}}\,{B}^{\mathrm{Prop}}.\,\left[A+B\right]\\\
{\mathrm{or}}_{\mathrm{introl}}&=\lambda{A}^{\mathrm{Prop}}\,{B}^{\mathrm{Prop}}\,{x}^{A}.\,\left\langle{\mathrm{inl}}\,{x}\right\rangle\\\
{\mathrm{or}}_{\mathrm{intror}}&=\lambda{A}^{\mathrm{Prop}}\,{B}^{\mathrm{Prop}}\,{y}^{B}.\,\left\langle{\mathrm{inr}}\,{y}\right\rangle\\\
\mathrm{Ex}&=\lambda{A}^{\mathrm{Type}}\,{P}^{A\rightarrow\mathrm{Prop}}.\,\left[\sum_{x{:}A}{P}\,{x}\right]\\\
{\mathrm{ex}}_{\mathrm{intro}}&=\lambda{A}^{\mathrm{Type}}\,{P}^{A\rightarrow\mathrm{Prop}}\,{x}^{A}\,{p}^{{P}\,{x}}.\,\left\langle{\left(x,p\right)}_{{\mbox{\scriptsize{$x${:}$A$.$P$}}}}\right\rangle\end{aligned}\right.$
Note how, because of the brackets, existentials and disjunctions are
prohibited from being eliminated to non-propositional types. Thanks to the
singleton rules, however, conjunction and falsity do not require brackets.
As a final example, consider the type Ascending n p of ascending sequences of
integers between p and n defined by mutual recursion with the proposition Ge m
p which stands from m is greater than or equal to p:
Inductive Ascending : Nat ${\rightarrow}$ Nat ${\rightarrow}$ Type := ${\mid}$
top : forall n, Ascending n n
${\mid}$ up : forall n p m, Ge m (s p) ${\rightarrow}$ Ascending n m
${\rightarrow}$ Ascending n p
with Ge : Nat ${\rightarrow}$ Nat ${\rightarrow}$ Prop := ${\mid}$ ascend :
forall m p, Ascending m p ${\rightarrow}$ Ge m p.
As Ascending has type Type, whereas Ge has type Prop, the translation to a
single inductive type is not as straightforward as Even and Odd. The
translation requires the use of brackets around the recursive calls:
$\left\\{\begin{aligned}
\mathrm{AscendingGe}&=\begin{array}[]{l}\mu{X}^{\left(\mathrm{Nat}\times\mathrm{Nat}\right)+\left(\mathrm{Nat}\times\mathrm{Nat}\right)\rightarrow\mathrm{Type}}.\,\lambda{i}^{\left(\mathrm{Nat}\times\mathrm{Nat}\right)+\left(\mathrm{Nat}\times\mathrm{Nat}\right)}.\,\mbox{\phantom{a}}\\\
\begin{array}[]{cl}&{\left({\sum}_{n{:}\mathrm{Nat}}^{{\mathrm{inl}}\,{\left(n,n\right)}}1\right)}\,{i}\\\
+&{\left({\sum}_{j{:}\mathrm{Nat}\times\mathrm{Nat}}^{{\mathrm{inl}}\,{j}}\sum_{m{:}\mathrm{Nat}}\left[{X}\,{\left({\mathrm{inr}}\,{\left(m,{s}\,{\left({{\pi}_{2}}\,{j}\right)}\right)}\right)}\right]\times{X}\,{\left({\mathrm{inl}}\,{\left({{\pi}_{1}}\,{j},m\right)}\right)}\right)}\,{i}\\\
+&{\left({\sum}_{j{:}\mathrm{Nat}\times\mathrm{Nat}}^{{\mathrm{inr}}\,{j}}{X}\,{\left({\mathrm{inl}}\,{j}\right)}\right)}\,{i}\\\
\end{array}\\\ \end{array}\\\
\mathrm{Ascending}&=\lambda{n}\,{p}.\,{\mathrm{AscendingGe}}\,{\left({\mathrm{inl}}\,{\left(n,p\right)}\right)}\\\
\mathrm{Ge}&=\lambda{m}\,{p}.\,\left[{\mathrm{AscendingGe}}\,{\left({\mathrm{inr}}\,{\left(m,p\right)}\right)}\right]\end{aligned}\right.$
## 5\. Impredicative Set
In addition to the impredicative sort $\mathrm{Prop}$, Coq has a sort
$\mathrm{Set}$ which is predicative by default but can be turned impredicative
with a flag. Where $\mathrm{Prop}$ is meant to be used in the context of
separating static and dynamic information, the spirit of the impredicative
sort $\mathrm{Set}$ is to be as powerful as possible without being
inconsistent. In the algebraic presentation, that means being stable by every
construction except dependent sums with the first projection in an arbitrary
sort (_strong sums_).
To mirror the optional nature of the impredicativity of $\mathrm{Set}$, the
rules for a predicative sort $\mathrm{Set}$ are given in Figure 6; to turn
impredicativity on, the rules of Figure 7 must be used _in addition_ to those
of predicative $\mathrm{Set}$. This presentation makes immediately apparent
that impredicative $\mathrm{Set}$ is an extension of predicative
$\mathrm{Set}$, in that every program of the latter typechecks in the former.
The rules of $\mathrm{Set}$ are the same as those of $\mathrm{Prop}$, with the
exception of $A+B$ which is in $\mathrm{Set}$ when both $A$ and $B$ are – even
with predicative $\mathrm{Set}$. Hence, there are types in $\mathrm{Set}$ with
several elements – _e.g._ $\mathbb{B}$. As a consequence, the bracketing
operation which coerces types in $\mathrm{Type}$ to $\mathrm{Set}$ does not
enjoy an explanation in terms of proof irrelevance, as was the case in
$\mathrm{Prop}$. As a matter of fact, there is no clear set-theoretical
description at all. A close cousin of $\mathrm{Set}$ bracketing, however, can
be found in homotopy type theory [15], where, roughly, groupoids are
_truncated_ to sets through a quotient of their homsets by the total relation.
$\Gamma\vdash\mathrm{Set}:\mathrm{Type}$
$\Gamma\vdash\prod_{x{:}A}B:\mathrm{Set}\lx@proof@logical@and\Gamma\vdash
A:\mathrm{Set}\Gamma,x{:}A\vdash B:\mathrm{Set}$
$\Gamma\vdash\sum_{x{:}A}P:\mathrm{Set}\lx@proof@logical@and\Gamma\vdash
A:\mathrm{Set}\Gamma,x{:}A\vdash P:\mathrm{Set}$
$\Gamma\vdash A+B:\mathrm{Set}\lx@proof@logical@and\Gamma\vdash
A:\mathrm{Set}\Gamma\vdash B:\mathrm{Set}$ $\Gamma\vdash 1:\mathrm{Set}$
$\Gamma\vdash 0:\mathrm{Set}$
$\Gamma\vdash\mu{X}^{A\rightarrow\mathrm{Set}}.\,F:A\rightarrow\mathrm{Set}\lx@proof@logical@and\Gamma\vdash
A:s\Gamma,X{:}A\rightarrow\mathrm{Set}\vdash
F:A\rightarrow\mathrm{Set}{{\mbox{{sp}}}_{X}}\,{F}$
$\Gamma\vdash{\sum}_{x{:}A}^{f}T:B\rightarrow\mathrm{Set}\lx@proof@logical@and\Gamma\vdash
A:\mathrm{Set}\Gamma\vdash B:\mathrm{Set}\Gamma,x{:}A\vdash
T:\mathrm{Set}\Gamma,x{:}A\vdash f:B$
$\Gamma\vdash\nu{X}^{A\rightarrow\mathrm{Set}}.\,F:A\rightarrow\mathrm{Set}\lx@proof@logical@and\Gamma\vdash
A:s\Gamma,X{:}A\rightarrow\mathrm{Set}\vdash
F:A\rightarrow\mathrm{Set}{{\mbox{{sp}}}_{X}}\,{F}$
$\Gamma\vdash{\left\\{A\right\\}}_{\mathrm{Set}}:\mathrm{Type}\Gamma\vdash
A:\mathrm{Set}$
$\Gamma\vdash{\mathrm{elt}}\,{u}:{\left\\{A\right\\}}_{\mathrm{Set}}\Gamma\vdash
u:A$
$\Gamma\vdash\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{elt}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}:{P}{\left[\,x\,{\setminus}\,u\,\right]}\lx@proof@logical@and\Gamma\vdash
u:{\left\\{A\right\\}}_{\mathrm{Set}}\Gamma,x{:}{\left\\{A\right\\}}_{\mathrm{Set}}\vdash
P:s\Gamma,y{:}A\vdash
v:{P}{\left[\,x\,{\setminus}\,{\mathrm{elt}}\,{y}\,\right]}$
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
{\mathrm{elt}}\,{u}\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\
x\leavevmode\nobreak\ \mathrm{return}\leavevmode\nobreak\
P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\mathrm{elt}}\,{y}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}\leadsto{v}{\left[\,y\,{\setminus}\,u\,\right]}$
${{\mbox{{sp}}}_{X}}\,{{\left\\{A\right\\}}_{\mathrm{Set}}}{{\mbox{{sp}}}_{X}}\,{A}$
Figure 6. Rules for predicative $\mathrm{Set}$
$\Gamma\vdash\prod_{x{:}A}B:\mathrm{Set}\lx@proof@logical@and\Gamma\vdash
A:s\Gamma,x{:}A\vdash B:\mathrm{Set}$
${{\mbox{{sp}}}_{X}}\,{{\left[A\right]}_{\mathrm{Set}}}{{\mbox{{sp}}}_{X}}\,{A}$
$\Gamma\vdash{\left[A\right]}_{\mathrm{Set}}:\mathrm{Set}\Gamma\vdash
A:\mathrm{Type}$ $\Gamma\vdash{\left\langle
u\right\rangle}_{\mathrm{Set}}:{\left[A\right]}_{\mathrm{Set}}\Gamma\vdash
u:A$
$\Gamma\vdash\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\
u\leavevmode\nobreak\ \mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\left\langle
y\right\rangle}_{\mathrm{Set}}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}:{P}{\left[\,x\,{\setminus}\,u\,\right]}\lx@proof@logical@and\Gamma\vdash
u:{\left[A\right]}_{\mathrm{Set}}\Gamma,x{:}{\left[A\right]}_{\mathrm{Set}}\vdash
P:\mathrm{Set}\Gamma,y{:}A\vdash v:{P}{\left[\,x\,{\setminus}\,{\left\langle
y\right\rangle}_{\mathrm{Set}}\,\right]}$
$\begin{array}[]{@{}l}\mathrm{match}\leavevmode\nobreak\ {\left\langle
u\right\rangle}_{\mathrm{Set}}\leavevmode\nobreak\
\mathrm{as}\leavevmode\nobreak\ x\leavevmode\nobreak\
\mathrm{return}\leavevmode\nobreak\ P\leavevmode\nobreak\ \mathrm{with}\\\
\begin{array}[]{|l@{~\Rightarrow~}l}{\left\langle
y\right\rangle}_{\mathrm{Set}}\hfil\leavevmode\nobreak\
\Rightarrow\leavevmode\nobreak\ &v\\\ \end{array}\\\
\end{array}\leadsto{v}{\left[\,y\,{\setminus}\,u\,\right]}$
Figure 7. Rules for impredicative $\mathrm{Set}$
## 6\. Conclusion
The algebraic presentation of Coq makes the conversion between sorts explicit.
The toplevel inductive definitions of Coq can be understood as implicitly
inserting canonical bracketing operations when an inductive type is declared
inside an impredicative sort but should be of a different sort due to its
form; and inserting type coercion from a smaller sort to a bigger sort when
applying a cumulativity rule.
Monolithic type definitions like in Coq have a number of advantages over the
algebraic presentation, they boil down to better type errors due to naming,
better type inference and better memory representation due to $n$-ary sums and
products. However, the value of the implicit coercions between sorts is less
clear. In particular, the bracketing operation to impredicative sorts is
probably a better guide for program extraction than the current method of
figuring whether or not a given type is a proposition, which interacts badly
with universe polymorphism [7]. Explicit coercions for extraction are also in
the spirit of [5].
All of the algebraic type constructors can actually be defined in Coq, except
the two fixed-points because there is no way to abstract over strictly
positive type families. So is it clear that expressions of the algebraic
presentation which do not use inductive or co-inductive fixed points can be
translated into Coq. Occurrences of the fixed points in a type must be
${\lambda}$-lifted and given a toplevel name. Some care must be given to
avoiding the duplication of such definitions otherwise types which must be
convertible for the expression to typecheck, might be seen as different in the
Coq translation. Apart from this technicality, translation from the algebraic
presentation to Coq is straightforward. We claim that, at least if we extend
the algebraic presentation to a hierarchy of universes and the strict
positivity condition is made a bit more fine-grained, Coq terms can be,
conversely, translated into the algebraic presentation.
## References
* [1] Andreas Abel. A polymorphic lambda-calculus with sized higher-order types. PhD thesis, 2006.
* [2] Steve Awodey and Andrej Bauer. Propositions as [Types]. Journal of Logic and Computation, 14(4):447–471, August 2004.
* [3] Henk Barendregt. Lambda calculus with types. Handbook of logic in computer science, 1992.
* [4] Bruno Barras. Semantical Investigations in Intuitionistic Set Theory and Type Theories with Inductive Families. Thèse d’Habilitation, 2013.
* [5] Bruno Barras and Bruno Bernardo. The implicit calculus of constructions as a programming language with dependent types. Foundations of Software Science and Computational Structures, 4962:365–379, 2008.
* [6] Gyesik Lee and Benjamin Werner. Proof-irrelevant model of CC with predicative induction and judgmental equality. 2011\.
* [7] Pierre Letouzey and Bas Spitters. Implicit and noncomputational arguments using monads. pages 1–16, 2005.
* [8] Peter Morris. Constructing Universes for Generic Programming. PhD thesis, 2007.
* [9] Peter Morris and Thorsten Altenkirch. Constructing strictly positive families. CATS ’07 Proceedings of the thirteenth Australasian symposium on Theory of computing, pages 111–121, 2007.
* [10] Peter Morris, Thorsten Altenkirch, and Neil Ghani. A universe of strictly positive families. IInternational journal of foundations of computer science, pages 83–107, 2009.
* [11] Ulf Norell. Towards a practical programming language based on dependent type theory. PhD thesis, Chalmers University of Technology, 2007.
* [12] Chris Okasaki. Purely functional data structures. 1999\.
* [13] Christine Paulin-Mohring. Définitions inductives en théorie des types d’ordre supérieur. PhD thesis, Université Claude Bernard-Lyon I, 1996.
* [14] The Coq development team. The Coq Proof Assistant.
* [15] The Univalent Foundations Program. Homotopy Type Theory: Univalent Foundations of Mathematics.
|
arxiv-papers
| 2013-09-23T11:14:12 |
2024-09-04T02:49:51.363899
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Hugo Herbelin and Arnaud Spiwack",
"submitter": "Arnaud Spiwack",
"url": "https://arxiv.org/abs/1309.5767"
}
|
1309.5853
|
11institutetext: Space Weather Laboratory, NASA Goddard Space Flight Center,
Greenbelt, MD, USA
11email: [email protected]; [email protected]
22institutetext: Max-Planck-Institut für Sonnensystemforschung, Max-Planck-
Strasse 2, 37191 Katlenburg-Lindau, Germany
22email: [email protected] 33institutetext: National Solar Observatory,
Sunspot, NM 88349, USA.
33email: [email protected]; [email protected]
# First use of synoptic vector magnetograms for global nonlinear force free
coronal magnetic field models
Tilaye Tadesse 11 T. Wiegelmann 22 S. Gosain 33 P. MacNeice 11 Alexei A.
Pevtsov 33
(Received ********** / Accepted ********)
###### Abstract
Context. The magnetic field permeating the solar atmosphere is generally
thought to provide the energy for much of the activity seen in the solar
corona, such as flares, coronal mass ejections (CMEs), etc. To overcome the
unavailability of coronal magnetic field measurements, photospheric magnetic
field vector data can be used to reconstruct the coronal field. Currently
there are several modelling techniques being used to calculate three-dimension
of the field lines into the solar atmosphere.
Aims. For the first time, synoptic maps of photospheric vector magnetic field
synthesized from Vector Spectromagnetograph (VSM) on Synoptic Optical Long-
term Investigations of the Sun (SOLIS) are used to model the coronal magnetic
field and estimate free magnetic energy in the global scale. The free energy
(i.e., the energy in excess of the potential field energy) is one of the main
indicators used in space weather forecasts to predict the eruptivity of active
regions.
Methods. We solve the nonlinear force-free field equations using optimization
principle in spherical geometry. The resulting three-dimensional magnetic
fields are used to estimate the magnetic free energy content
$E_{\mathrm{free}}=E_{\mathrm{nlfff}}-E_{\mathrm{pot}}$, i.e., the difference
of the magnetic energies between the nonpotential field and the potential
field in the global solar corona. For comparison, we overlay the extrapolated
magnetic field lines with the extreme ultraviolet (EUV) observations by the
Atmospheric Imaging Assembly on board SDO.
Results. For a single Carrington rotation 2121, we find that the global NLFFF
magnetic energy density is $10.3\%$ higher than the potential one. Most of
this free energy is located in active regions.
Conclusions.
###### Key Words.:
Magnetic fields – Sun: corona – Sun: photosphere – methods: numerical
## 1 Introduction
In the solar coronal plasma, magnetic energy is the prime energy reservoir
that fuels the dynamical evolution of eruptive events, but it remains an open
question how the magnetic energy is released. The amount of energy associated
with the magnetic field is much larger than other energy sources, and the
dynamics of the coronal configuration is determined by the evolution of its
magnetic field (Forbes 2000; Low 2001).
Free energy is defined as excess of energy as compared with potential field.
In fact, one can in principle release more energy than free energy by, for
example, annihilating magnetic field. Also, numerically, one can get pre-post
flare difference in energy larger than free magnetic energy if magnetic field
at the photosphere changes (e.g., some flux elements disappear). It quantifies
the energy deviation of the coronal magnetic field from its potential state
(Metcalf et al. 2005; Régnier 2007; Aschwanden 2012). The magnetic free energy
is stored in the form of electric currents flowing along the magnetic field.
Free magnetic energy of solar solar magnetic fields can be affected by several
processes such as, e.g., photospheric shearing flows, magnetic flux emergence
and magnetic reconnection (e.g., Welsch 2006; Fang et al. 2012).
To understand the role that the magnetic field plays in energizing the solar
corona, it is important to calculate the amount of free energy in order to
quantify the energy budget in a catastrophic energy release event, as well as
for estimating upper limits in forecasting individual events in real-time.
Using various extrapolation techniques for the coronal magnetic field under
the assumption of force-free fields, the spatial and temporal evolution of the
coronal magnetic free energy during solar flares has been extensively studied
(Régnier & Canfield 2006; Guo et al. 2008; Jing et al. 2010; Tadesse et al.
2012a; Meyer et al. 2013). From the 3D coronal magnetic configurations, we can
derive the magnetic energy in the corona as:
$E_{\mathrm{M}}=\frac{1}{8\pi}\int_{V}\textbf{B}\cdot\textbf{B}\,\,r^{2}\sin\theta
drd\theta d\phi$ (1)
The free magnetic energy in spherical geometry is calculated by computing the
nonpotential field $B_{\mathrm{nlfff}}(r,\theta,\phi)$ with a numerical
nonlinear force-free field (NLFFF) code and a potential field
$B_{\mathrm{pot}}(r,\theta,\phi)$ for the same photospheric boundary data, so
that the difference of the magnetic field energy density in the coronal volume
V encompassing the active regions of interest can be quantified as
$E_{\mathrm{M_{free}}}=E_{\mathrm{M_{nlfff}}}-E_{\mathrm{M_{pot}}}$.
Since the corona is optically thin, direct measurements of the 3-D magnetic
field are very difficult to implement and interpret. Therefore, the present
observations for the magnetic fields based on the spectropolarimetric method
(the Zeeman and the Hanle effects) are limited to low layers of solar
atmosphere (photosphere and chromosphere). Even if direct measurement
techniques for the 3-D magnetic fields in the chromosphere and the corona have
considerably improved in recent decades (Lin et al. 2000, 2004; Liu & Lin
2000), further developments are needed before accurate data are routinely
available. Thus, the problem of measuring the coronal field and its embedded
electrical currents leads us to use numerical modelling to infer the field
strength in the higher layers of the solar atmosphere from the measured
photospheric field.
Force-free extrapolation of photospheric magnetic fields is currently used as
the primary tool for the modeling of coronal magnetic fields. In this model
assumption, the corona magnetic forces are dominant so that all non-magnetic
forces like plasma pressure gradient and gravity can be neglected in the
lowest order (Gary 2001). This implies that, if appreciable currents are
present, these must be aligned with the magnetic field, since otherwise the
resulting Lorentz forces could not be balanced by the nonmagnetic forces. The
equilibrium structure of the static coronal magnetic field can be described
using the force-free assumption as:
$(\nabla\times\textbf{B})=\alpha\textbf{B}\Rightarrow(\nabla\times\textbf{B})\times\textbf{B}=0$
(2) $\nabla\cdot\textbf{B}=0$ (3)
where B is the magnetic field. The force-free parameter $\alpha$ of Equation
(2) can be a function of position, but the combination of Equation (2) and (3)
($\textbf{B}\cdot\nabla\alpha=0$) requires that $\alpha$ be constant along a
given field line. Potential ($\alpha=0$) and linear force-free fields
(whenever $\alpha$ is constant everywhere in the volume under consideration)
can be used as a first step to model the general structure of magnetic fields
in the solar corona. Practically the pre-eruptive magnetic fields are
nonlinear force-free fields ($\alpha$ being a function of position) as
supported by both observational and theoretical reasons. For details of those
models we direct the readers to Wiegelmann & Sakurai (2012).
Nonlinear force-free field codes have been routinely applied to the
reconstruction of the coronal field of a single active region using the
Cartesian geometry. In that case, the curvature of the solar surface does not
play an important role. Solar Dynamics Observatory (SDO) mission has made
repeated observations of large scale events in which connections between
widely separated active regions may play fundamental role (Martens et al.
2012). Even before SDO, it was known that large-scale connectivity is
important for solar eruptive and non-eruptive activity (e.g., studies of
sympathetic flares, transequatorial loops (Pevtsov 2000), effects of distant
active regions on large-scale coronal brightness (Pevtsov & Acton 2001),
eruption of filaments triggered by remote flux emergence and evolution (e.g.,
Balasubramaniam et al. 2011). Therefore, this needs motivate us to implement a
NLFFF procedure in spherical geometry (Wiegelmann 2007; Tadesse et al. 2009,
2011, 2012a, 2012b; Guo et al. 2012; Tadesse et al. 2013a; Amari et al. 2013)
In this study, we estimate the free magnetic energy for global corona using
data from SOLIS/VSM. We compare the extrapolated potential and NLFFF magnetic
loops with extreme ultraviolet (EUV) observations by the Atmospheric Imaging
Assembly (AIA) on board SDO. This comparison helps to identify whether the
NLFFF model reconstructs the magnetic configuration better than the potential
field model in the global scale. In this paper, we present some descriptions
of the dataset used for analysis in Section 2. The spherical optimization
procedure used for modeling 3-D magnetic field in global corona is presented
in Section 3. Then, we present results of our studies in Section 4. A summary
and discussions are finally presented in Section 5.
## 2 Instrumentation and data set
In this paper, we use the first ever synoptic vector magnetograms, i.e.,
Carrington maps of the three components of the magnetic field vector, the
radial $B_{r}$, the poloidal $B_{\theta}$, and the toroidal $B_{\phi}$, to
perform first global nonlinear force free field extrapolation based on
optimization algorithm (Wiegelmann 2007; Tadesse et al. 2009). The synoptic
vector field maps are derived using daily full disk photospheric vector
magnetograms from Vector Spectromagnetograph (VSM) instrument of Synoptic
Optical Long-term Investigations of the Sun (SOLIS), a synoptic observing
facility (Keller et al. 2003; Balasubramaniam & Pevtsov 2011). Further detail
about these synoptic maps and their properties can be found in Gosain et al.
(2013). Here we briefly describe the instrument characteristics and the full
disk vector field observations.
VSM routinely obtains full disk magnetograms in photospheric and chromospheric
lines as a part of the synoptic program of National Solar Observatory i.e.,
NSO Integrated Synoptic Program (NISP). In order to obtain full disk
photospheric vector magnetograms SOLIS/VSM measures Stokes $(S=I,Q,U,V)$
profiles in photospheric Fe I $630.15-630.25$ nm line pair. A single full disk
scan (2048 scan lines) takes only about 20 minutes, thanks to the long slit of
the spectrograph which intersects the solar disk from one limb to another in
one shot. The spatial sampling is 1 arc-sec per pixel with square pixels. The
spectral sampling is 2.4pm per pixel. A single Stokes set $(I,Q,U,V)$ per slit
position is obtained typically in about 0.6 second. The telescope itself is
designed to be free of instrumental polarization by employing symmetric
optical configuration and performing polarization modulation just after the
prime focus, after the slit. Dual beam analysis of polarization using
polarizing beam splitter avoids seeing induced cross talks in the signal. The
polarimeter calibration is done routinely to calibrate the Stokes vector for
cross-talks. The signal to noise ratio (SNR) in the continuum of the Stokes
profiles is typically $>1000$. The magnetic field vector is inferred from the
calibrated Stokes profiles by performing inversion in the framework of Milne-
Eddington model for stellar atmosphere following Unno-Rachkovsky formalism
(Skumanich & Lites 1987).
Only pixels with polarization signal above the threshold of $0.1\%$ of
continuum intensity, Ic, are inverted to obtain the magnetic (field strength,
inclination angle, and azimuth angle) and thermodynamic (e.g., Doppler width,
Doppler velocity, source function, temperature) parameters. The threshold of
$0.1\%$ of Ic corresponds to typical noise level in the continuum. Using this
threshold avoids fitting profiles buried in the noise.
Further details about instrument and pipeline reduction steps can be found
elsewhere (e.g., Jones et al. 2002; Henney et al. 2006; Balasubramaniam &
Pevtsov 2011). The noise in SOLIS magnetograms is estimated to be few Gauss in
longitudinal and 70 G in the transverse field measurements (Tadesse et al.
2013a). The 180 degree azimuth ambiguity is resolved using a different
(faster) ambiguity resolution method developed by Rudenko & Anfinogentov
(2011).
Figure 1: Synoptic Carrington maps of the vector magnetic field components are
shown for CR-2121. The panels from top to bottom show the distribution of the
$B(r)$, $B(\phi)$ and $B(\theta)$ components, respectively. The $B_{r}$ map is
scaled between $\pm$ 100 G, and the $B_{\phi}$ and $B_{\theta}$ maps are
scaled to $\pm$ 20 G. The positive values of $B_{r}$, $B_{\phi}$ and
$B_{\theta}$ point, respectively, upward, to the right (westward) and
southward.
## 3 Magnetic field modeling
For modelling the coronal magnetic field in a global scale, we use the
variational principle originally proposed by Wheatland et al. (2000).
Wiegelmann (2007) has developed an optimization method to reconstruct the
NLFFF for global solar corona by minimizing an objective functional $L$ that
combines Lorentz forces and the divergence of the magnetic field in spherical
geometry. The code has been tested with semi-analytic force-free solutions
(Low & Lou 1990). If the functional is minimized to zero, Equations (2) and
(3) are satisfied simultaneously. The optimization procedure in the spherical
geometry has been implemented by Tadesse et al. (2009), for restricted area
with large field of views. Later Tadesse et al. (2011) modified the objective
functional of the optimization method (Wiegelmann & Inhester 2010) for
spherical geometry as
$L=L_{f}+L_{d}+\nu L_{\mathrm{photo}}$ (4)
$L_{\mathrm{f}}=\int_{V}B^{-2}\big{|}(\nabla\times{\textbf{B}})\times{\textbf{B}}\big{|}^{2}r^{2}\sin\theta
drd\theta d\phi$
$L_{\mathrm{d}}=\int_{V}\big{|}\nabla\cdot{\textbf{B}}\big{|}^{2}r^{2}\sin\theta
drd\theta d\phi$
$L_{\mathrm{photo}}=\int_{S}\big{(}\textbf{B}-\textbf{B}_{\mathrm{obs}}\big{)}\cdot\textbf{W}(\theta,\phi)\cdot\big{(}\textbf{B}-\textbf{B}_{\mathrm{obs}}\big{)}r^{2}\sin\theta
d\theta d\phi$
where $L_{f}$ and $L_{d}$ measure how well the force-free Eqs. (2) and
divergence-free (3) conditions are fulfilled, respectively. The main reason
for modification of the code was that we need to deal with boundary data of
different noise levels and qualities or even miss some data points completely.
Hence, the third integral, $L_{\mathrm{photo}}$, is the surface integral over
the photosphere which allows us to relax the field on the photosphere towards
force-free solution without too much deviation from the original surface field
data.
Figure 2: Global field lines of (a) the potential field model and (b) the
NLFFF model overlaid on the AIA 171Å image. Green and red lines represent open
and closed magnetic field lines, respectively.
SOLIS/VSM provides full disk vector magnetograms from which Synoptic
Carrington maps of the vector magnetic field components are synthesized.
However, for pixels below the polarization threshold the inversion was not
performed and field data there will be missing for these pixels (see Figure
1). Typically, the field is missing where its magnitude is small; thus these
pixels would have a small impact on the model even if they were measured
correctly. Within the error margin of a measured field value, any value is
just as good as any other, and from this range of values we take the value
that fits the force-free field best. In order to treat those pixels with
missing data, we used the diagonal matrix, $\textbf{W}(\theta,\phi)$, which
gives different weights to the observed surface field components depending on
the relative accuracy in measurement. In this sense, missing data is
considered most inaccurate and is taken into account by setting
$\textbf{W}(\theta,\phi)$ to zero in all elements of the matrix.
Photospheric magnetic field has a plasma-$\beta$ of order unity, which does
not satisfy the force-free condition (Gary 2001). Therefore, the vector
magnetogram data are inconsistent with the force-free assumption, which is
absolutely essential condition for NLFFF extrapolation. To find suitable
boundary conditions for the NLFFF field modeling, we have to preprocess the
measured synoptic vector magnetograms by using a preprocessing scheme
developed by Tadesse et al. (2009) in spherical geometry. This preprocessing
scheme removes forces and torques from the boundary and approximates the
photospheric magnetic field to the low plasma-$\beta$ force-free chromosphere.
For a detailed description of the current code implementation, we refer to
Wiegelmann (2007) and Tadesse et al. (2011).
## 4 Results
This study requires extrapolating the three-dimensional potential and NLFFF
coronal fields from the photospheric boundary in global scale. We use synoptic
maps of photospheric vector magnetic field observed during 4-31 March, 2012.
During this observation there were about 25 active regions all over the solar
globe. In order to use our spherical optimization code for global corona, we
adopt a uniform spherical grid $r$, $\theta$, $\phi$ with $n_{r}=300$,
$n_{\theta}=450$, and $n_{\phi}=900$ grid points in the direction of radius,
latitude, and longitude, respectively, with the field of view of
$[r_{\rm{min}}=1R_{\odot}:r_{\rm{max}}=2.5R_{\odot}]\times[\theta_{\rm{min}}=-85^{\circ}:\theta_{\rm{max}}=85^{\circ}]\times[\phi_{\rm{min}}=0^{\circ}:\phi_{\rm{max}}=360^{\circ}]$.
The code solves the NLFFF equations in the bounded domain between $1R_{\odot}$
and the source surface at $2.5R_{\odot}$. The outer boundary is kept fixed
using the initial potential field values. All current-carrying field lines
have to close inside the volume. The domain outside $2.5R_{\odot}$ is not
included in the model, because the force-free approach is no longer justified
here. The magnetic field extrapolation has been carried out almost for global
corona by excluding the polar regions where the magnetic field measurements
are highly being influenced by noise. Before performing NLFFF extrapolations,
we use the preprocessed radial magnetic field component ${\bf{B}}_{r}$ to
compute the corresponding potential field using a spherical harmonic expansion
for initializing our spherical NLFFF code during relaxation towards a force-
free state in the computational volume.
Figure 3: Field lines of (a) the potential field model and (b) the NLFFF model
around ARs 11429 and 11430 overlaid on the AIA 171Å image. Green and red lines
represent open and closed magnetic field lines, respectively.
Figure 4: The magnetic field line skeletons (a) of the entire solar globe from
the NLFFF model and image of the sun observed by SOHO/LASCO C2 coronagraph at
16:33UT.
The main purpose of this work is to study the structures of the global
potential and NLFFF magnetic fields and to estimate free magnetic energy
available to power solar eruptions during Carrington rotation 2121. In
addition, we compare which of those two models is best agree with observation
in global environment. To do this, we plot the selected fieldlines of the
potential and NLFFF models in Figure 2. We overlay the field lines with an AIA
171 Å image. The field lines of the potential and NLFFF models are
reconstructed from the same footpoints. The potential field lines in Figure
3(a) have an obvious deviation from the observed EUV loops, since the
projection of the field lines indicated by blue and black arrows divate from
EUV loops. However, those loops are best overlaid by NLFFF lines than
potential ones (see Figure 3(b)). Therefore, the qualitative comparison
between the model magnetic field lines and the observed EUV loops indicates
that the NLFFF model provides a more consistent field for global corona
magnetic field reconstruction. Figure 4 shows that there is over all
similarity between the corresponding NLFFF model field lines and image of the
sun observed by SOHO/LASCO C2 coronagraph. However, one can see that NLFFF
does not represent well a linear structure above coronal helmet in low-left
corner in Figure 3(b). This could be due to the missing data from the polar
region.
Figure 5: Evolution of the entire functional L (blue line) and its three term
in Eq. 4 during the optimization process. The black line corresponds to
$L_{\rm photo}$, the red line to $L_{\rm f}$, and the green line to $L_{\rm
d}$.
In addition to the above comparison to quantify the degree of disagreement
between the two model vector field solutions in the global corona volume that
are specified on the identical sets of grid points, we use the vector
correlation metric ($C_{\rm vec}$) which is also used analogous to the
standard correlation coefficient for scalar functions. The correlation was
calculated (Schrijver et al. 2006) from
$C_{\mathrm{vec}}=\frac{\sum_{i}\textbf{v}_{i}\cdot\textbf{u}_{i}}{\Big{(}\sum_{i}|\textbf{v}_{i}|^{2}\sum_{i}|\textbf{u}_{i}|^{2}\Big{)}^{1/2}},$
(5)
where $\textbf{v}_{i}$ and $\textbf{u}_{i}$ are the vectors at each grid point
$i$. If the vector fields are identical, then $C_{\rm vec}=1$; if
$\textbf{v}_{i}\perp\textbf{u}_{i}$ , then $C_{\rm vec}=0$. The degree of
convergence towards a force-free and divergence-free model solution can be
quantified by the integral measures of the Lorentz force and divergence terms
in the minimization functional in Equation (4), computed over the entire solar
globe. $L_{\rm f}$ and $L_{\rm d}$ of Equation (4) measure how well the force-
free and divergence-free conditions are fulfilled, respectively. In Table 1,
we provide some quantitative measures to rate the quality of our
reconstruction. Column $1$ names the corresponding models. Columns $2-3$ show
how well the force-balance and solenoidal conditions are fulfilled for both
models. Figure 5 shows how well the functional $L$ converge to zero during
iteration process. In the last column, the vector correlation shows that there
is disagreement between the two model field solutions.
The energy stored in the magnetic field as a result of field line stressing
into a nonpotential configuration has been identified as the source of flare
energy. Therefore, to understand the physics of solar flares, including the
local reorganisation of the magnetic field and the acceleration of energetic
particles, we have first to estimate the free magnetic energy available for
such phenomena. This free magnetic energy can be converted into kinetic and
thermal energy. We estimate the free magnetic energy, the difference between
the extrapolated NLFFF and the potential field with the same normal boundary
conditions in the photosphere. We therefore estimate the upper limit to the
free magnetic energy associated with coronal currents of the form
$E_{\mathrm{free}}=\frac{1}{8\pi}\int_{V}\Big{(}{B^{2}}_{\mbox{nlff}}-{B^{2}}_{\mbox{pot}}\Big{)}dV\,\,$
(6)
${B}_{\mbox{pot}}$ and ${B}_{\mbox{nlff}}$ represent the potential and NLFFF
magnetic field, respectively. Our result for the estimation of free-magnetic
energy in Table 2 shows that the NLFFF model has $10.3\%$ more energy than the
corresponding potential field model. Figure 6 shows iso-surface plot of free
magnetic energy density in the volume above synoptic map. There are strong
free energy concentrations above each active region over the solar globe.
Figure 6: Iso-surfaces (ISs) of free magnetic energy density ($8.5\times 10^{20}$ erg) computed within the entire computational domain. Table 1: Evaluation of the reconstruction quality for the potential field and NLFFF models. We have used spherical grids of $300\times 450\times 900$. Model | $L_{\rm f}$ | $L_{\rm d}$ | $L_{\rm photo}$ | $C_{\rm vec}$
---|---|---|---|---
Potential | $0.000$ | $0.000$ | $0.001$ | $1$
NLFFF | $0.391$ | $0.697$ | $0.302$ | $0.893$
Table 2: The magnetic energy associated with the 3-D potential and NLFFF field configurations calculated from synoptic vector magnetogram. Model | $E_{\rm{total}}(10^{33}{\rm{erg}})$ | $E_{\rm{free}}(10^{33}{\rm{erg}})$
---|---|---
Potential | $41.6$ | $0$
NLFFF | $46.4$ | $4.8$
## 5 Conclusion and outlook
Most of the NLFFF procedures are implemented in the Cartesian coordinates.
Therefore, both potential and nonlinear force-free field (NLFFF) codes in
Cartesian geometry are not well suited for larger domains, since the spherical
nature of the solar surface cannot be neglected when the field of view is
large. Therefore, it is necessary to implement a NLFFF procedure in spherical
geometry for use when large-scale boundary data are in use.
In this study, we have investigated the coronal magnetic field and free
magnetic energy associated with global corona by analyzing Carrington synoptic
maps of photospheric vector magnetic field synthesized from Vector
Spectromagnetograph (VSM) on Synoptic Optical Long-term Investigations of the
Sun (SOLIS) which has been observed during 4-31 March, 2012. The Carrington
rotation number for this observation is 2121. During this particular
observation, there were about ten active regions distributed across the globe.
We have used our spherical NLFFF and potential codes to compute the magnetic
field solutions over global corona. This is the first NLFFF magnetic field
extrapolation for the global corona using real data.
We have compared the magnetic field solutions from both potential and NLFFF
models. The qualitative comparison between the model magnetic field lines and
the observed EUV loops indicates that the NLFFF model provides a more
consistent field for global corona magnetic field reconstruction. For this
particular Carrington rotation we find that the global NLFFF magnetic energy
density is $10.3\%$ higher than the potential one. For future, we have a plan
to study the evolution of global free magnetic energy. In this study, most of
this free energy is located in active regions.
###### Acknowledgements.
This work utilizes SOLIS data obtained by the NSO Integrated Synoptic Program
(NISP), managed by the National Solar Observatory, which is operated by the
Association of Universities for Research in Astronomy(AURA), Inc. under a
cooperative agreement with the National Science Foundation. This research was
supported by an appointment to the NASA Postdoctoral Program at the Goddard
Space Flight Center (GSFC), administered by Oak Ridge Associated Universities
through a contract with NASA. The work of T. Wiegelmann was supported by DLR-
grant $50$ OC $453$ $0501$.
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|
arxiv-papers
| 2013-09-23T15:40:40 |
2024-09-04T02:49:51.375187
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Tilaye Tadesse and T. Wiegelmann and S. Gosain and P. MacNeice and\n Alexei A. Pevtsov",
"submitter": "Tilaye Tadesse",
"url": "https://arxiv.org/abs/1309.5853"
}
|
1309.5886
|
# Generalized three-intensity decoy state method for measurement device
independent quantum key distribution
Zong-Wen Yu1,2, Yi-Heng Zhou1, and Xiang-Bin Wang1,3111Email Address:
[email protected] 1State Key Laboratory of Low Dimensional Quantum
Physics, Tsinghua University, Beijing 100084, People s Republic of China
2Data Communication Science and Technology Research Institute, Beijing 100191,
China
3 Shandong Academy of Information and Communication Technology, Jinan 250101,
People s Republic of China
###### Abstract
We present improved results and more general results for the decoy-state
measurement device independent quantum key distribution (MDI-QKD) where the
intensities of all 3 pulses can be non-zero. We present a more tightened
explicit formula for the lower bound of the yield of two single-photon pulses
for this generalized 3-intensity protocol, which can be applied to the
recently proposed MDI-QKD with imperfect single-photon sources such as the
coherent states or the heralded states from the parametric down-conversion. We
strictly prove that our result is better than the prior art result. With this
result, the final key rate rises drastically.
###### pacs:
03.67.Dd, 42.81.Gs, 03.67.Hk
## I Introduction
Recently, much attention has been spent on the security of real set-ups in
quantum key distribution (QKD) BB84 ; GRTZ02 . After the proposal of the decoy
state method ILM ; H03 ; wang05 ; LMC05 ; AYKI ; haya ; peng ; wangyang ; rep
; njp , the security of QKD with an imperfect single photon source is
guaranteed given whatever Eve’s attack to the channel, including the famous
photon-number splitting attack PNS1 ; PNS .
However, the guaranteed security assumes that Eve cannot attack device inside
Alice or Bob’s lab. As was demonstrated lyderson , given a lossy channel and
limited detection efficiency, Eve can attack the gated-mode detectors inside
Alice or Bob’s lab if there is no additional device to detect the abnormally
strong light at Alice or Bob’s lab. To reach the goal of ultimate security,
one can, in principle use the device independent protocol ind1 . However,
these protocols seem to be technically demanding at the moment.
Recently, measurement device independent QKD (MDI-QKD) was proposed based on
the idea of entanglement swapping ind3 ; ind2 . If we want to obtain a higher
key rate, we can choose to directly use an imperfect single-photon source such
as the coherent state ind2 with decoy-state method for this, say the MDI
decoy-state method. Calculation formulas for the practical decoy-state
implementation with only a few different states has been studied in, e.g.,
Refs. wangPRA2013 ; qing3 , and then further studied both experimentally
tittel1 ; tittel2 ; liuyang and theoretically lopa ; han ; curty ; Wang2013 .
Theoretical study for the decoy-state MDI-QKD with only a few intensities are
of particularly interesting. The first formula for 3-intensity decoy-state
MDI-QKD was proposed in wangPRA2013 . It assumes that each user uses three
different states, one vacuum state and two non-vacuum states. The results can
apply to sources such as the coherent states, the parametric down conversion,
and so on. It was then extended to the case that the 3 intensities are all
non-vacuum coherent states lopa . One can see that the major formula there is
identical with the one in Ref.wangPRA2013 in the case when Alice and Bob use
coherent pulses and the weakest pulse is vacuum. In Ref.Wang2013 ; Wang201309
, we presented an improved explicit formulas of 3-state decoy-state method for
the MDI-QKD with the weakest pulse being vacuum.
Here in this work, we shall present an improved results for the decoy-state
MDI-QKD where all 3 intensities are non-vacuum. Compared with the prior-art
result lopa , our improved formulas here make a more tightened faithful
estimation of the yield and the phase flip rate of the single-photon pulse
pairs and hence can raise the final key rate drastically. Also, our result is
for a general type of sources with a simple condition, it does not limit to
the coherent-state source. In the second section, we present the improved
explicit formulas for 3-intensity decoy-state MDI-QKD and prove analytically
that our result is better than the previous lopa . In the third section, we
give the numerical simulations. The article is ended with a concluding remark.
## II Generalized Three-intensity decoy state method for MDI-QKD
In the protocol, each time a pulse-pair (two-pulse state) is sent to the relay
for detection. The relay is controlled by an un-trusted third party (UTP). The
UTP will announce whether the pulse-pair has caused a successful event. Those
bits corresponding to successful events will be post-selected and further
processed for the final key. Since real set-ups only use imperfect single-
photon sources, we need the decoy-state method for security.
We assume Alice (Bob) has three sources, $v_{A},d_{A},s_{A}$
($v_{B},d_{B},s_{B}$) which can only emit three different states
$\rho_{v_{A}},\rho_{d_{A}},\rho_{s_{A}}$
($\rho_{v_{B}},\rho_{d_{B}},\rho_{s_{B}}$), respectively, in photon number
space. It’s worth pointing out that our protocol is applicable to the weakest
pulse being vacuum or not. Suppose
$\left\\{\begin{array}[]{lcr}\rho_{v_{A}}=\sum_{k}a^{v}_{k}|k\rangle\langle
k|,&&\rho_{v_{B}}=\sum_{k}b^{v}_{k}|k\rangle\langle k|,\\\
\rho_{d_{A}}=\sum_{k}a^{d}_{k}|k\rangle\langle
k|,&&\rho_{d_{B}}=\sum_{k}b^{d}_{k}|k\rangle\langle k|,\\\
\rho_{s_{A}}=\sum_{k}a^{s}_{k}|k\rangle\langle
k|,&&\rho_{s_{B}}=\sum_{k}b^{s}_{k}|k\rangle\langle k|.\end{array}\right.$ (1)
At each time, Alice will randomly select one of her 3 sources to emit a pulse,
and so does Bob. The pulse form Alice and the pulse from Bob form a pulse pair
and are sent to the UTP. We regard equivalently that each time a two-pulse
source is selected and a pulse pair (one pulse from Alice, one pulse from Bob)
is emitted. There are many different two-pulse sources used in the protocol.
We denote $\alpha\beta$ for the two pulse source when the pulse-pair is
produced by source $\alpha$ at Alice’s side and source $\beta$ at Bob’s side,
$\alpha$ can be one of $\\{v_{A},d_{A},s_{A}\\}$ and $\beta$ can be one of
$\\{v_{B},d_{B},s_{B}\\}$. For example, at a certain time $j$ Alice uses
source $v_{A}$ and Bob uses source $s_{B}$, we say the pulse pair is emitted
by source $v_{A}s_{B}$.
In the protocol, two different bases (e.g., $Z$ basis consisting of horizontal
polarization $|H\rangle\langle H|$ and vertical polarization $|V\rangle\langle
V|$, and $X$ basis consisting of $\pi/4$ and $3\pi/4$ polarizations) are used.
The density operator in photon number space alone does not describe the state
in the composite space. We shall apply the decoy-state method analysis in the
same basis (e.g., $Z$ basis or $X$ basis) for pulses from sources
$\alpha\beta(\alpha\in\\{v_{A},d_{A},s_{A}\\},\beta\in\\{v_{B},d_{B},s_{B}\\})$.
Therefore we only need consider the density operators in the photon number
space.
According to the decoy-state theory, the yield of a certain set of pulse pairs
is defined as the happening rate of a successful event (announced by the UTP)
corresponding to pulse pairs out of the set. We use capital letter
$Y_{\alpha\beta}$ to denote the nine yields for sources
$\alpha\beta(\alpha\in\\{v_{A},d_{A},s_{A}\\},\beta\in\\{v_{B},d_{B},s_{B}\\})$
and these values can be directly calculated from the observed experimental
data. We use $y_{mn}$ to denotes the yield given that Alice and Bob send
respectively an $m$-photon and $n$-photon pulses.
Our first major task is to deduce $y_{11}$ from the known values, i.e., to
formulate $y_{11}$, the yield of state $|1\rangle\langle
1|\otimes|1\rangle\langle 1|$ in the nine known values
$\\{Y_{\alpha\beta}|\alpha\in\\{v_{A},d_{A},s_{A}\\},\beta\in\\{v_{B},d_{B},s_{B}\\}\\}$.
Using the convex proposition, we have the following relations
$Y_{\alpha\beta}=\sum_{k,l\geq 0}a^{\alpha}_{k}b^{\beta}_{l}y_{kl},$ (2)
where $\alpha\in\\{v_{A},d_{A},s_{A}\\}$ and
$\beta\in\\{v_{B},d_{B},s_{B}\\}$. In practice, the imperfect single-photon
sources have the property that the probability of emitting an empty pulse is
greater than the probability of the no-empty pulse. With this fact, in order
to estimate the lower bound of $y_{11}$, we need to eliminate the unknown
values $y_{0l}(l\geq 0)$ and $y_{k0}(k\geq 0)$ firstly.
Without causing any ambiguity, we omit the subscripts A and B in the following
of this paper. Cancel out the terms $y_{0k}$ and $y_{j0}$ using Gaussian
elimination, we obtain the following four relations
$\displaystyle\left\\{\begin{array}[]{c}a_{0}^{v}b_{0}^{v}\hat{Y}_{dd}-a_{0}^{v}b_{0}^{d}\hat{Y}_{dv}-a_{0}^{d}b_{0}^{v}\hat{Y}_{vd}+a_{0}^{d}b_{0}^{d}\hat{Y}_{vv}=0,\\\
a_{0}^{v}b_{0}^{v}\hat{Y}_{ds}-a_{0}^{v}b_{0}^{s}\hat{Y}_{dv}-a_{0}^{d}b_{0}^{v}\hat{Y}_{vs}+a_{0}^{d}b_{0}^{s}\hat{Y}_{vv}=0,\\\
a_{0}^{v}b_{0}^{v}\hat{Y}_{sd}-a_{0}^{v}b_{0}^{d}\hat{Y}_{sv}-a_{0}^{s}b_{0}^{v}\hat{Y}_{vd}+a_{0}^{s}b_{0}^{d}\hat{Y}_{vv}=0,\\\
a_{0}^{v}b_{0}^{v}\hat{Y}_{ss}-a_{0}^{v}b_{0}^{s}\hat{Y}_{sv}-a_{0}^{s}b_{0}^{v}\hat{Y}_{vs}+a_{0}^{s}b_{0}^{s}\hat{Y}_{vv}=0,\end{array}\right.$
(7)
where $\hat{Y}_{ij}=Y_{ij}-\sum_{k,l\geq
1}a_{k}^{i}b_{l}^{j}y_{kl},(i,j\in\\{v,d,s\\})$. For the sake of convenience,
we denote
$\displaystyle\left\\{\begin{array}[]{c}\tilde{Y}_{dd}^{\prime}=a_{0}^{v}b_{0}^{v}{Y}_{dd}-a_{0}^{v}b_{0}^{d}{Y}_{dv}-a_{0}^{d}b_{0}^{v}{Y}_{vd}+a_{0}^{d}b_{0}^{d}{Y}_{vv},\\\
\tilde{Y}_{ds}^{\prime}=a_{0}^{v}b_{0}^{v}{Y}_{ds}-a_{0}^{v}b_{0}^{s}{Y}_{dv}-a_{0}^{d}b_{0}^{v}{Y}_{vs}+a_{0}^{d}b_{0}^{s}{Y}_{vv},\\\
\tilde{Y}_{sd}^{\prime}=a_{0}^{v}b_{0}^{v}{Y}_{sd}-a_{0}^{v}b_{0}^{d}{Y}_{sv}-a_{0}^{s}b_{0}^{v}{Y}_{vd}+a_{0}^{s}b_{0}^{d}{Y}_{vv},\\\
\tilde{Y}_{ss}^{\prime}=a_{0}^{v}b_{0}^{v}{Y}_{ss}-a_{0}^{v}b_{0}^{s}{Y}_{sv}-a_{0}^{s}b_{0}^{v}{Y}_{vs}+a_{0}^{s}b_{0}^{s}{Y}_{vv},\end{array}\right.$
(12)
and introduce the following notations
$h_{x_{k}}^{i}=\frac{x_{k}^{i}}{x_{0}^{i}},\quad(k\geq
0;x\in\\{a,b\\};i\in\\{v,d,s\\}).$ (13)
Divide the four relations in Eq.(7) by the positive factors
$a_{0}^{v}a_{0}^{d}b_{0}^{v}b_{0}^{d}$,
$a_{0}^{v}a_{0}^{d}b_{0}^{v}b_{0}^{s}$, $a_{0}^{v}a_{0}^{s}b_{0}^{v}b_{0}^{d}$
and $a_{0}^{v}a_{0}^{s}b_{0}^{v}b_{0}^{s}$ respectively, we obtain the
following simplified equation
$\displaystyle\left\\{\begin{array}[]{c}\tilde{Y}_{dd}=\sum_{k,l\geq
1}(h_{a_{k}}^{d}-h_{a_{k}}^{v})(h_{b_{l}}^{d}-h_{b_{l}}^{v})y_{kl},\\\
\tilde{Y}_{ds}=\sum_{k,l\geq
1}(h_{a_{k}}^{d}-h_{a_{k}}^{v})(h_{b_{l}}^{s}-h_{b_{l}}^{v})y_{kl},\\\
\tilde{Y}_{sd}=\sum_{k,l\geq
1}(h_{a_{k}}^{s}-h_{a_{k}}^{v})(h_{b_{l}}^{d}-h_{b_{l}}^{v})y_{kl},\\\
\tilde{Y}_{ss}=\sum_{k,l\geq
1}(h_{a_{k}}^{s}-h_{a_{k}}^{v})(h_{b_{l}}^{s}-h_{b_{l}}^{v})y_{kl}.\end{array}\right.$
(18)
where $\tilde{Y}_{ij}$ are defined as follows
$\tilde{Y}_{ij}=\frac{\tilde{Y}_{ij}^{\prime}}{a_{0}^{v}a_{0}^{i}b_{0}^{v}b_{0}^{j}},\quad(i,j\in\\{d,s\\}),$
(19)
with $\tilde{Y}_{ij}^{\prime}$ being presented in Eq.(12). Furthermore, if we
introduce
$\tilde{h}_{x_{k}}^{i}=h_{x_{k}}^{i}-h_{x_{k}}^{v},\quad(k\geq
0;x\in\\{a,b\\};i\in\\{v,d,s\\}),$ (20)
with $h_{x_{k}}^{i}$ defined in Eq.(13), we can write the relations about
variables $y_{kl}$ given by Eq.(18) into the following simplified form
$\displaystyle\tilde{Y}_{dd}$ $\displaystyle=$
$\displaystyle\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{d}y_{11}+\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{2}}^{d}y_{12}+\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}y_{21}+J_{dd},$
(21) $\displaystyle\tilde{Y}_{ds}$ $\displaystyle=$
$\displaystyle\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{s}y_{11}+\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{2}}^{s}y_{12}+\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{s}y_{21}+J_{ds},$
(22) $\displaystyle\tilde{Y}_{sd}$ $\displaystyle=$
$\displaystyle\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{1}}^{d}y_{11}+\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{2}}^{d}y_{12}+\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{d}y_{21}+J_{sd},$
(23) $\displaystyle\tilde{Y}_{ss}$ $\displaystyle=$
$\displaystyle\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{1}}^{s}y_{11}+\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{2}}^{s}y_{12}+\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{s}y_{21}+J_{ss},$
(24)
where $\tilde{Y}_{ij}$ is defined in Eq.(12), $\tilde{h}_{x_{k}}^{i}$ is given
by Eq.(20), and $J_{ij}=\sum_{(k,l)\in
J_{1}}{\tilde{h}_{a_{k}}^{i}\tilde{h}_{b_{l}}^{j}y_{kl}},(i,j\in\\{d,s\\})$
with $J_{1}=\\{(k,l)|k,l\geq 1;k+l\geq 4\\}$.
With these preparations, now we are going to introduce the following very
important condition:
$\frac{\tilde{h}_{x_{k}}^{s}}{\tilde{h}_{x_{k}}^{d}}\geq\frac{\tilde{h}_{x_{2}}^{s}}{\tilde{h}_{x_{2}}^{d}}\geq\frac{\tilde{h}_{x_{1}}^{s}}{\tilde{h}_{x_{1}}^{d}},\quad(x\in\\{a,b\\})$
(25)
for $k\geq 2$. We can easily prove that the imperfect sources used in practice
such as the coherent state source, the heralded source out of the parametric-
down conversion, satisfy the above restriction.
In order to get a lower bound of $y_{11}$, we should derive the expression of
$y_{11}$ with Eqs.(21-24) firstly. Combining Eqs.(21-23), we obtain the
expression of $y_{11}$ by eliminating $y_{12}$ and $y_{21}$ such that
$y_{11}=y_{11}^{(123)}+\sum_{(m,n)\in J_{1}}f_{11}^{(123)}(m,n)y_{mn},$ (26)
where $J_{1}=\\{(m,n)|m\geq 1,n\geq 1,m+n\geq 4\\}$,
$y_{11}^{(123)}=\frac{(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})\tilde{Y}_{dd}-\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})\tilde{Y}_{ds}-\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{d}(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})\tilde{Y}_{sd}}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})},$
(27)
and
$f_{11}^{(123)}(m,n)=\frac{\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{n}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{m}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{m}}^{d})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})+\tilde{h}_{a_{m}}^{d}\tilde{h}_{b_{1}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})(\tilde{h}_{b_{2}}^{d}\tilde{h}_{b_{n}}^{s}-\tilde{h}_{b_{2}}^{s}\tilde{h}_{b_{n}}^{d})}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})}.$
(28)
In these expressions, we use the superscript $*^{(123)}$ to denote the result
obtained with the first three equations from Eqs.(21-24). Under the conditions
presented in Eq.(25), we can easily find out that
$(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})\geq
0$,
$(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})\geq
0$,
$(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{m}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{m}}^{d})\geq
0$ for all $m\geq 1$ and
$(\tilde{h}_{b_{2}}^{d}\tilde{h}_{b_{n}}^{s}-\tilde{h}_{b_{2}}^{s}\tilde{h}_{b_{n}}^{d})\geq
0$ for all $n\geq 2$. Then we know that $f_{11}^{(123)}(m,n)\geq 0$ hold for
all $(m,n)\in J_{1}$. With this fact, we obtain a lower bound from Eq.(26) by
setting $y_{mn}=0,(m,n)\in J_{1}$ such that
$\underline{y_{11}}=y_{11}^{(123)}\leq y_{11},$ (29)
where $y_{11}^{(123)}$ is defined by Eq.(27). This and Eq.(27) are our major
formulas for the decoy-state method implementation for MDI-QKD in this
section.
Similarly, we can get other expressions with choosing any other three
equations from Eqs.(21-24). For example, we choose Eqs.(21-22,24). By
eliminating $y_{12}$ and $y_{21}$, we get another expression of $y_{11}$ such
that
$y_{11}=y_{11}^{(124)}+\sum_{(m,n)\in J_{1}}f_{11}^{(124)}(m,n)s_{mn},$ (30)
where
$y_{11}^{(124)}=\frac{\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{s}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})\tilde{Y}_{dd}+(\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})\tilde{Y}_{ds}-\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{d}(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})\tilde{Y}_{ss}}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{s}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})},$
(31)
and
$f_{11}^{(124)}(m,n)=\frac{\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{n}}^{s}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{m}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{m}}^{d})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})+\tilde{h}_{a_{m}}^{d}\tilde{h}_{b_{1}}^{s}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})(\tilde{h}_{b_{2}}^{d}\tilde{h}_{b_{n}}^{s}-\tilde{h}_{b_{2}}^{s}\tilde{h}_{b_{n}}^{d})}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{s}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})}.$
(32)
Under the conditions presented in Eq.(25), we can also find out that
$f_{11}^{(124)}(m,n)\geq 0$ for all $(m,n)\in J_{1}$. Then we know that
$y_{11}^{(124)}$ is alos a lower bound of $y_{11}$. On the other hand, by
comparing $f_{11}^{(123)}(m,n)$ and $f_{11}^{(124)}(m,n)$, we have
$\displaystyle f_{11}^{(123)}(m,n)-f_{11}^{(124)}(m,n)$ $\displaystyle=$
$\displaystyle\frac{\tilde{h}_{a_{2}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{m}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{m}}^{d})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{n}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{n}}^{d})}{-\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{1}}^{s}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})}\leq
0,$
for all $(m,n)\in J_{1}$. Then we know that
$y_{11}^{(123)}\geq y_{11}^{(124)},$ (33)
with Eq.(26) and Eq.(30). With the relation presented in Eq.(33) we know that
the lower bound $y_{11}^{(123)}$ is tighter than the lower bound
$y_{11}^{(124)}$. In the same way, we can get another two lower bounds
$y_{11}^{(134)}$ and $y_{11}^{(234)}$ of $y_{11}$ with Eqs.(21,23-24) and
Eqs.(22-24) respectively. Furthermore, we can also prove that
$y_{11}^{(123)}\geq y_{11}^{(134)},\quad y_{11}^{(123)}\geq y_{11}^{(234)}.$
(34)
Now we only consider Eq.(21) and Eq.(24). By eliminating $y_{12}$ or $y_{21}$
respectively, we get two expressions of $y_{11}$ such that
$\displaystyle y_{11}$ $\displaystyle=$ $\displaystyle
y_{11}^{(14a)}+\sum_{(m,n)\in J_{2}}f_{11}^{(14a)}(m,n)y_{mn},$ (35)
$\displaystyle y_{11}$ $\displaystyle=$ $\displaystyle
y_{11}^{(14b)}+\sum_{(m,n)\in J_{2}}f_{11}^{(14a)}(m,n)y_{mn},$ (36)
where $J_{2}=\\{(m,n)|m\geq 1,n\geq 1,m+n\geq 3\\}$,
$\displaystyle y_{11}^{(14a)}$ $\displaystyle=$
$\displaystyle\frac{\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{2}}^{s}\tilde{Y}_{dd}-\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{2}}^{d}\tilde{Y}_{ss}}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{1}}^{s}(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})},$
(37) $\displaystyle y_{11}^{(14b)}$ $\displaystyle=$
$\displaystyle\frac{\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{s}\tilde{Y}_{dd}-\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}\tilde{Y}_{ss}}{\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{1}}^{s}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})},$
(38)
and
$\displaystyle f_{11}^{(14a)}(m,n)$ $\displaystyle=$
$\displaystyle\frac{\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{2}}^{d}\tilde{h}_{a_{m}}^{s}\tilde{h}_{b_{n}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{2}}^{s}\tilde{h}_{a_{m}}^{d}\tilde{h}_{b_{n}}^{d}}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{1}}^{s}(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})},$
(39) $\displaystyle f_{11}^{(14b)}(m,n)$ $\displaystyle=$
$\displaystyle\frac{\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}\tilde{h}_{a_{m}}^{s}\tilde{h}_{b_{n}}^{s}-\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{s}\tilde{h}_{a_{m}}^{d}\tilde{h}_{b_{n}}^{d}}{\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{1}}^{s}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})}.$
(40)
For any sources used in the protocol, we must have either
$K_{a}=\frac{\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{2}}^{s}}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{2}}^{d}}\leq\frac{\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{s}}{\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}}=K_{b}\quad\textrm{or}\quad
K_{a}>K_{b}.$
Suppose the former one holds, we can easily find out that
$f_{11}^{(14a)}(m,n)\geq 0$ for all $(m,n)\in J_{2}$ and $y_{11}^{(14a)}$ is a
lower bound of $y_{11}$. On the other hand, if $K_{a}\geq K_{b}$ holds, we
have $f_{11}^{(14b)}(m,n)\geq 0$ for all $(m,n)\in J_{2}$ and $y_{11}^{(14b)}$
is a lower bound of $y_{11}$. Considering the following two relations
$K_{a}-K_{b}=\frac{\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d}}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{d}},$
(41)
and
$\displaystyle f_{11}^{(14a)}(m,n)-f_{11}^{(14b)}(m,n)$ $\displaystyle=$
$\displaystyle\frac{\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{2}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{m}}^{s}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{n}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{m}}^{d}\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{n}}^{d})(K_{a}-K_{b})}{\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{1}}^{s}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})},$
we know that $K_{a}-K_{b}$ and $f_{11}^{(14a)}-f_{11}^{(14b)}$ have the same
sign which means that they are both positive or negative simultaneously. Then
we can write the lower bound of $y_{11}$ with Eq.(21) and Eq.(24) into the
following compact form
$y_{11}^{(14)}=\min\\{y_{11}^{(14a)},y_{11}^{(14b)}\\}.$ (42)
If both Alice and Bob use coherent pulses, the lower bound given by Eq.(42) is
just the result presented in Refs. lopa ; curty . In the coming, we will prove
that the lower bound $y_{11}^{(123)}$ given in Eq.(27) is more tightly than
$y_{11}^{(14)}$. Firstly, if we suppose $K_{a}\leq K_{b}$ holds, then we know
that
$\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}\leq\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d}$
and $y_{11}^{(14)}=y_{11}^{(14a)}$. For any $(m,n)\in J_{1}$ we have
$\displaystyle f_{11}^{(123)}(m,n)-f_{11}^{(14a)}(m,n)$ $\displaystyle=$
$\displaystyle-\frac{(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{m}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{m}})D_{a}}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{1}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})},$
where
$D_{a}=(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{d}\tilde{h}_{b_{n}}^{s}+\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d}\tilde{h}_{b_{n}}^{d}+\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}\tilde{h}_{b_{n}}^{d}+\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{d}\tilde{h}_{b_{n}}^{s})\geq\tilde{h}_{b_{2}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{2}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{2}}^{d})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{n}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{n}}^{d})$.
Then we know that
$\displaystyle f_{11}^{(123)}(m,n)-f_{11}^{(14a)}(m,n)$ $\displaystyle\leq$
$\displaystyle-\frac{\tilde{h}_{b_{2}}^{d}(\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{m}}^{s}-\tilde{h}_{a_{1}}^{s}\tilde{h}_{a_{m}}^{d})(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{n}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{n}}^{d})}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{1}}^{d}(\tilde{h}_{b_{1}}^{d}\tilde{h}_{b_{2}}^{s}-\tilde{h}_{b_{1}}^{s}\tilde{h}_{b_{2}}^{d})}\leq
0.$
We can easily know that $y_{11}^{(123)}\geq y_{11}^{(14)}$ when $K_{a}\leq
K_{b}$ with this equation. Secondly, if we suppose $K_{a}>K_{b}$ holds, we can
easily prove that $f_{11}^{(123)}(m,n)-f_{11}^{(14b)}(m,n)\leq 0$ for all
$(m,n)\in J_{1}$ within the same way. Then we get $y_{11}^{(123)}\geq
y_{11}^{(14)}$ when $K_{a}>K_{b}$. This completes the proof that
$y_{11}^{(123)}\geq y_{11}^{(14)}$.
In order to estimate the final key rate, we also need the upper bound of error
rate caused by the two single-photon pulses, say $e_{11}$. Similar to the
total gain, the total error rate with source $\alpha\beta$ chosen by Alice and
Bob can be written as ind2
$\displaystyle\tilde{T}_{dd}$ $\displaystyle=$
$\displaystyle\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{d}t_{11}+\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{2}}^{d}t_{12}+\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{d}t_{21}+K_{dd},$
(43) $\displaystyle\tilde{T}_{ds}$ $\displaystyle=$
$\displaystyle\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{s}t_{11}+\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{2}}^{s}t_{12}+\tilde{h}_{a_{2}}^{d}\tilde{h}_{b_{1}}^{s}t_{21}+K_{ds},$
(44) $\displaystyle\tilde{T}_{sd}$ $\displaystyle=$
$\displaystyle\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{1}}^{d}t_{11}+\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{2}}^{d}t_{12}+\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{d}t_{21}+K_{sd},$
(45) $\displaystyle\tilde{T}_{ss}$ $\displaystyle=$
$\displaystyle\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{1}}^{s}t_{11}+\tilde{h}_{a_{1}}^{s}\tilde{h}_{b_{2}}^{s}t_{12}+\tilde{h}_{a_{2}}^{s}\tilde{h}_{b_{1}}^{s}t_{21}+K_{ss},$
(46)
where $T_{ij}=Y_{ij}E_{ij},(i,j\in\\{v,d,s\\})$, $t_{mn}=y_{mn}e_{mn}$, and
$\displaystyle\tilde{T}_{dd}$ $\displaystyle=$
$\displaystyle\frac{T_{dd}}{a_{0}^{d}b_{0}^{d}}-\frac{T_{dv}}{a_{0}^{d}b_{0}^{v}}-\frac{T_{vd}}{a_{0}^{v}b_{0}^{d}}+\frac{T_{vv}}{a_{0}^{v}b_{0}^{v}},$
(47) $\displaystyle\tilde{T}_{ds}$ $\displaystyle=$
$\displaystyle\frac{T_{ds}}{a_{0}^{d}b_{0}^{s}}-\frac{T_{dv}}{a_{0}^{d}b_{0}^{v}}-\frac{T_{vs}}{a_{0}^{v}b_{0}^{s}}+\frac{T_{vv}}{a_{0}^{v}b_{0}^{v}},$
(48) $\displaystyle\tilde{T}_{sd}$ $\displaystyle=$
$\displaystyle\frac{T_{sd}}{a_{0}^{s}b_{0}^{d}}-\frac{T_{sv}}{a_{0}^{s}b_{0}^{v}}-\frac{T_{vd}}{a_{0}^{v}b_{0}^{d}}+\frac{T_{vv}}{a_{0}^{v}b_{0}^{v}},$
(49) $\displaystyle\tilde{T}_{ss}$ $\displaystyle=$
$\displaystyle\frac{T_{ss}}{a_{0}^{s}b_{0}^{s}}-\frac{T_{sv}}{a_{0}^{s}b_{0}^{v}}-\frac{T_{vs}}{a_{0}^{v}b_{0}^{s}}+\frac{T_{vv}}{a_{0}^{v}b_{0}^{v}},$
(50)
and $K_{ij}=\sum_{(k,l)\in
J_{1}}{\tilde{h}_{a_{k}}^{i}\tilde{h}_{b_{l}}^{j}t_{kl}},(i,j\in\\{d,s\\})$
with $J_{1}=\\{(k,l)|k,l\geq 1,k+l\geq 4\\}$. According to Eq.(43), we can
find out the upper bound of $e_{11}$ such that
$e_{11}\leq
e_{11}^{(1)}=\frac{\tilde{T}_{dd}}{\tilde{h}_{a_{1}}^{d}\tilde{h}_{b_{1}}^{d}y_{11}}=\overline{e_{11}}.$
(51)
In the protocol, there are two different bases. We denote $y_{11}^{Z}$ and
$y_{11}^{X}$ for yields of single-photon pulse pairs in the $Z$ and $X$ bases,
respectively. Consider those post-selected bits cased by source $d_{A}d_{B}$
in the $Z$ basis. After an error test, we know the bit-flip error rate of this
set, say $T_{d_{A}d_{B}}^{Z}=Y_{d_{A}d_{B}}^{Z}E_{d_{A}d_{B}}^{Z}$. We also
need the phase-flip rate for the subset of bits which are caused by the two
single-photon pulse, say $e_{11}^{ph}$, which is equal to the flip rate of
post-selected bits caused by a single photon in the $X$ basis, say
$e_{11}^{X}$. Given this, we can now calculate the key rate by the well-know
formula. For example, for those post-selected bits caused by source
$s_{A}s_{B}$, it is
$R=a_{1}^{s}b_{1}^{s}y_{11}^{Z}[1-H(e_{11}^{X})]-Y_{s_{A}s_{B}}^{Z}fH(E_{s_{A}s_{B}}^{Z}),$
(52)
where $f$ is the efficiency factor of the error correction method used.
## III Numerical Simulation
Figure 1: (Color online) The relative value between estimated parameter of
$y_{11}$ and the asymptotic limit of the infinite decoy-state method versus
the total channel transmission loss using 3-intensity decoy state MDI-QKD. In
the simulations, we assume that $\mu_{v}=\nu_{v}=0.01$, $\mu_{d}=\nu_{d}=0.1$,
$\mu_{s}=\nu_{s}=0.5$ and the value of other parameters are presented in Table
1.
Figure 2: (Color online) The estimated parameter of $e_{11}$ versus the total
channel transmission loss using 3-intensity decoy state MDI-QKD. In the
simulations, we assume that $\mu_{v}=\nu_{v}=0.01$, $\mu_{d}=\nu_{d}=0.1$,
$\mu_{s}=\nu_{s}=0.5$ and the value of other parameters are presented in Table
1.
Figure 3: (Color online) The estimated key rate $R$ versus channel
transmission using 3-intensity decoy state MDI-QKD. In the simulations, we
assume that $\mu_{v}=\nu_{v}=0.01$, $\mu_{d}=\nu_{d}=0.1$,
$\mu_{s}=\nu_{s}=0.5$ and the value of other parameters are presented in Table
1.
Figure 4: (Color online) The relative value between the optimal key rate
obtained with different methods and the asymptotic limit of the infinite
decoy-state method versus the total channel transmission loss using
3-intensity decoy state MDI-QKD. In the simulations, we assume that
$\mu_{v}=\nu_{v}=0.01$, $\mu_{d}=\nu_{d}=0.1$, $\mu_{s}=\nu_{s}$ and the value
of other parameters are presented in Table 1. The optimal key rate is just the
value with maximizing the key rate with $\mu_{s}=\nu_{s}\in(\mu_{d},1)$.
Figure 5: (Color online) The optimal intensities versus the total channel transmission loss using 3-intensity decoy state MDI-QKD. In the simulations, we assume that $\mu_{v}=\nu_{v}=0.01$, $\mu_{d}=\nu_{d}=0.1$, $\mu_{s}=\nu_{s}$ and the value of other parameters are presented in Table 1. The optimal densities of the signal-state pulses are obtained by maximizing the key rate. Table 1: List of experimental parameters used in numerical simulations: $e_{0}$ is the error rate of background, $e_{d}$ is the misalignment-error probability; $p_{d}$ is the dark count rate per detector; $f$ is the error correction inefficiency. $e_{0}$ | $e_{d}$ | $p_{d}$ | $f$
---|---|---|---
0.5 | 1.5% | $3.0\times 10^{-6}$ | 1.16
Now, we present some numerical simulations to comparing our results with the
existing results lopa ; curty . Below for simplicity, we suppose that Alice
and Bob use the coherent-state sources. Here, we denote Alice’s sources
$\\{v_{A},d_{A},s_{A}\\}$ by their intensities $\\{\mu_{v},\mu_{d},\mu_{s}\\}$
and Bob’s sources $\\{v_{B},d_{B},s_{B}\\}$ by their intensities
$\\{\nu_{v},\nu_{d},\nu_{s}\\}$ respectively. The UTP locates in the middle of
Alice and Bob, and the UTP’s detectors are identical, i.e., they have the same
dark count rate and detection efficiency, and their detection efficiency does
not depend on the incoming signals. We shall estimate what values would be
probably observed for the gains and error rates in the normal cases by the
linear models as in ind2 ; Wang2013 ; Wang201309 :
$\displaystyle|n\rangle\langle
n|=\sum_{k=0}^{n}C_{n}^{k}\xi^{k}(1-\xi)^{n-k}|k\rangle\langle k|$
where $\xi^{k}$ is the transmittance for a distance from Alice to the UTB. For
fair comparison, we use the same parameter values used in ind2 ; Wang2013 ;
Wang201309 for our numerical evaluation, which follow the experiment reported
in UrsinNP2007 . For simplicity, we shall put the detection efficiency to the
overall transmittance $\eta=\xi^{2}\zeta$. We assume all detectors have the
same detection efficiency $\zeta$ and dark count rate $p_{d}$. The values of
these parameters are presented in Table 1. With this, the total gains
$Y_{\mu_{i},\nu_{j}}^{\omega},(\omega=X,Z)$ and error rates
$Y_{\mu_{i},\nu_{j}}^{\omega}E_{\mu_{i},\nu_{j}}^{\omega},(\omega=X,Z)$ of
Alice’s intensity $\mu_{i}(i\in\\{v,d,s\\})$ and Bob’s intensity
$\nu_{j}(j\in\\{v,d,s\\})$ can be calculated. By using these values, we can
estimate the lower bounds of yield $y_{11}^{Z}$ with Eq.(29) and Eq.(42).
Also, we can estimate the upper bounds of error rate $e_{11}^{X}$ with Eq.(51)
where the lower bound of $y_{11}^{X}$ being estimated by Eq.(29) and Eq.(42).
The estimated parameters of $y_{11}$ and $e_{11}$ are shown in Fig.1 and
Fig.2, respectively, which clearly shows that our methods are more tightly
than the pre-existed results. In order to see more clearly, in Fig.1, we plot
the relative value of $y_{11}$ to the result obtained with the infinite decoy-
state method. We can observe that our results are more close to the asymptotic
limit of the infinite decoy-state method than the pre-existed results lopa ;
curty . Furthermore, with these parameters, we can estimate the final key rate
$R$ of this protocol with Eq.(52) which is shown in Fig.3. In these figures,
the blue dashed line is obtained by Eq.(42) which is just the results
presented in Refs. lopa ; curty , the red solid line is obtained by the method
presented in section 2 with Eq.(29), the black dash-dot line is the result
obtained by the infinite decoy-state method. In the simulation, the densities
used by Alice and Bob are assigned to $\mu_{v}=\nu_{v}=0.01$,
$\mu_{d}=\nu_{d}=0.1$, $\mu_{s}=\nu_{s}=0.5$.
Furthermore, if we fix the densities of the decoy-state pulses used by Alice
and Bob, the final key rate will change with Alice and Bob taking different
densities for their signal-state pulses. Then we can find out the optimal
densities of their signal-state pulses with maximizing the final key rate.
Here, we also take $\mu_{v}=\nu_{v}=0.01$, $\mu_{d}=\nu_{d}=0.1$ and assume
that $\mu_{s}=\nu_{s}\in(\mu_{d},1)$. In Fig.4, we plot the relative value of
the optimal key rate to the result obtained with the infinite decoy-state
method. We can observe that our result is better than the pre-existed results.
The optimal densities with the optimal key rate versus the total channel
transmission loss is given in Fig.5.
## IV Conclusion
We study the MDI-QKD in practice where the intensities of all 3 different
states can be nonzero. Our result here is the most general result for the
3-intensity method: Setting the weakest state to be zero, we obtain the result
in the special case where a vacuum state and two non-vacuum states are
usedWang2013 . The result here is not limited to the coherent states. It
applies to any states that satisfy Eq.(25), e.g., the heralded states from the
parametric down conversion. Our result is most efficient: It offers a more
tightened bound for $y_{11}$ and therefore a higher key rate than the prior
artlopa , as has been strictly proven and also numerically demonstrated. The
key rate can be further improved when it is combined with the method as shown
in Ref.Wang201309 .
Acknowledgement: We acknowledge the support from the 10000-Plan of Shandong
province, the National High-Tech Program of China Grants No. 2011AA010800 and
No. 2011AA010803 and NSFC Grants No. 11174177 and No. 60725416.
## References
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* (14) B. Huttner, N. Imoto, N. Gisin, and T. Mor, Phys. Rev. A 51, 1863 (1995); H.P. Yuen, Quantum Semiclassic. Opt. 8, 939 (1996).
* (15) L. Lyderson et al, Nature Photonics, 4, 686(2010); I. Gerhardt et al, Nature Commu. 2, 349 (2011).
* (16) D. Mayers and A. C.-C. Yao, in Proceedings of the 39th Annual Symposium on Foundations of Computer Science (FOCS98) (IEEE Computer Society, Washington, DC, 1998), p. 503; A. Acin et al., Phys. Rev. Lett. 98, 230501 (2007); Scarani V and Renner R Phys. Rev. Lett. 100, 302008 (2008); Scarani V and Renner R 2008 3rd Workshop on Theory of Quantum Computation, Communication and Cryptography (TQC 2008), (University of Tokyo, Tokyo 30 Jan C1 Feb 2008) See also arXiv:0806.0120
* (17) S.L. Braunstein and S. Pirandola, Phys. Rev. Lett. 108, 130502 (2012).
* (18) H.-K. Lo, M. Curty, and B. Qi, Phys. Rev. Lett., 108,130503(2012), K. Tamaki et al, Phys. Rev. A, 85, 042307 (2012).
* (19) Qin-Wang abd Xiang-Bin Wang, arXiv:1305.6480.
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* (21) Allison Rubenok, Joshua A. Slater, Philip Chan, Itzel Lucio-Martinez, Wolfgang Tittel, 1304.2463v1.
* (22) P. Chan, J. A. Slater, I. Lucio-Martinez, A. Rubenok, W. Tittel, arxiv:1204.0738v1.
* (23) Y. Liu et al, arXiv:1209.6178v1.
* (24) Feihu Xu, Marcos Curty, Bing Qi, Hoi-Kwong Lo , arXiv:1305.6965v1.
* (25) Chun Zhou, Wan-Su Bao, Wei Chen, Hong-Wei Li, Zhen-Qiang Yin, Yang Wang, Zheng-Fu Han, arXiv: 1308.3374v1.
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* (27) Z.-W. Yu, Y.-H. Zhou, X.-B. Wang, arXiv:1308.5677.
* (28) Z.-W. Yu, Y.-H. Zhou, X.-B. Wang, arXiv:1309.0471.
* (29) R. Ursin, F. Tiefenbacher, T. Schmitt-Manderbach, H. Weier, T. Scheidl, M. Lindenthal, B. Blauensteiner, T. Jennewein, J. Perdigues, P.Trojek, B. Oemer, M. F uerst, M. Meyenburg, J. Rarity, Z. Sodnik, C. Barbieri, H. Weinfurther, and A. Zeilinger, Nat. Phys. 3, 481 (2007).
|
arxiv-papers
| 2013-09-23T17:24:50 |
2024-09-04T02:49:51.384117
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Z.-W. Yu, Y.-H. Zhou, and Xiang-Bin Wang",
"submitter": "Xiang-Bin Wang",
"url": "https://arxiv.org/abs/1309.5886"
}
|
1309.5983
|
# Remarks on the tensor degree of finite groups
Ahmad M.A. Alghamdi Department of Mathematical Sciences, Faculty of Applied
Sciences
Umm Alqura University, P.O. Box 14035, Makkah, 21955, Saudi Arabia
[email protected] and Francesco G. Russo DEIM, Universitá degli Studi di
Palermo
Viale Delle Scienze, Edificio 8, 90128, Palermo, Italy
[email protected]
###### Abstract.
The present paper is a note on the tensor degree of finite groups, introduced
recently in literature. This numerical invariant generalizes the commutativity
degree through the notion of nonabelian tensor square. We show two
inequalities, which correlate the tensor and the commutativity degree of
finite groups, and, indirectly, structural properties will be discussed.
###### Key words and phrases:
Tensor degree, commutativity degree, exterior degree
###### 2010 Mathematics Subject Classification:
Primary: 20J99, 20D15; Secondary: 20D60; 20C25
## 1\. The relative tensor degree
All the groups of the present paper are supposed to be finite. Having in mind
the exponential notation for the conjugation of two elements $x$ and $y$ in a
group $G$, that is, the notation $x^{y}=y^{-1}xy$, we may follow [3, 4, 13] in
saying that two normal subgroups $H$ and $K$ of $G$ act compatibly upon each
other, if
$\left(h_{2}^{k_{1}}\right)^{h_{1}}=\left(\left({h_{2}}^{h_{1}^{-1}}\right)^{k_{1}}\right)^{h_{1}}\
\mathrm{and}\
\left(k_{2}^{h_{1}}\right)^{k_{1}}=\left(\left({k_{2}}^{k_{1}^{-1}}\right)^{h_{1}}\right)^{k_{1}}$
for all $h_{1},h_{2}\in H$ and $k_{1},k_{2}\in K$, and if $H$ and $K$ act upon
themselves by conjugation. Given $h\in H$ and $k\in K$, the nonabelian tensor
product $H\otimes K$ is the group generated by the symbols $h\otimes k$
satisfying the relations $h_{1}h_{2}\otimes
k_{1}=(h_{2}^{h_{1}}\otimes{k_{1}}^{h_{1}})\ (h_{1}\otimes k_{1})$ and
$k_{1}k_{2}\otimes h_{1}=(k_{1}\otimes h_{1})\ (h_{1}^{k_{1}}\otimes
k_{2}^{k_{1}})$ for all $h_{1},h_{2}\in H$ and $k_{1},k_{2}\in K$. The map
$\kappa_{H,K}:h\otimes k\in H\otimes K\mapsto[h,k]\in[H,K]$
turns out to be an epimorphism, whose kernel $\ker\kappa_{H,K}=J(G,H,K)$ is
central in $H\otimes K$. The reader may find more details and a topological
approach to $J(G,H,K)$ in [4, 5, 11, 13]. The short exact sequence
$\begin{CD}1@>{}>{}>J(G,H,K)@>{}>{}>H\otimes
K@>{\kappa_{H,K}}>{}>[H,K]@>{}>{}>1\end{CD}$
is a central extension. In the special case $G=H=K$, we have that
$J(G)=J(G,G,G)=\ker\kappa_{G,G}=\ker\kappa$ and $H\otimes K=G\otimes G$ is
called nonabelian tensor square of $G$. The fundamental properties of
$G\otimes G$ have been described in the classical paper [3], in which it is
noted that $\kappa:x\otimes y\in G\otimes G\mapsto\kappa(x\otimes y)=[x,y]\in
G^{\prime}$ is an epimorphism of groups with $\ker\kappa=J(G)$ and
$1\rightarrow J(G)\rightarrow G\otimes
G{\overset{\kappa}{\rightarrow}}G^{\prime}\rightarrow 1$ is a central
extension. The group $J(G)$ is important from the perspective of the algebraic
topology, in fact $J(G)\cong\pi_{3}(SK(G,1))$ is the third homotopy group of
the suspension of an Eilenberg–MacLane space $K(G,1)$ (see [4] for more
details).
As done in [12], we may consider the tensor centralizer
$C^{\otimes}_{K}(H)=\\{k\in K\ |\ h\otimes k=1,\ \forall h\in
H\\}=\bigcap_{h\in H}C^{\otimes}_{K}(h)$
and the tensor center $C^{\otimes}_{G}(G)=Z^{\otimes}(G)={\underset{x\in
G}{\bigcap}}C_{G}^{\otimes}(x)$ and one can check that $C^{\otimes}_{G}(x)$
and $Z^{\otimes}(G)$ are subgroups of $G$ such that
$C^{\otimes}_{G}(x)\subseteq C_{G}(x)$ and $Z^{\otimes}(G)\subseteq Z(G)$.
Generalizing what has been done in [12], we may define the relative tensor
degree
$d^{\otimes}(H,K)=\frac{|\\{(h,k)\in H\times K\ |\ h\otimes
k=1\\}|}{|H||K|}=\frac{1}{|H|\ |K|}\sum_{h\in H}|C^{\otimes}_{K}(h)|$
of $H$ and $K$. Notice that $d^{\otimes}(G)=1$ if and only if
$Z^{\otimes}(G)=G$. Unfortunately, few results are available on the relative
tensor degree at the moment and these are contained mainly in [12], where it
is discussed the tensor degree $d^{\otimes}(G)=d^{\otimes}(G,G)$ of $G$. On
the other hand, there is a rich literature (see for instance [1, 2, 7, 8, 9])
on the relative commutativity degree
$d(H,K)=\frac{|\\{(h,k)\in H\times K\ |\ [h,k]=1\\}|}{|H||K|}=\frac{1}{|H|\
|K|}\sum_{h\in H}|C_{K}(h)|=\frac{k_{K}(H)}{|H|}$
of $H$ and $K$ (not necessarily normal this time) of $G$. Here $k_{K}(H)$ is
the number of $K$–conjugacy classes that constitute $H$. In particular, if
$G=H=K$, we find the well known commutativity degree
$d(G)=d(G,G)=k_{G}(G)/|G|$. Our fist main contribution is the following.
###### Theorem 1.1.
Let $H,K$ be two normal subgroups of a group $G$ and $p$ the smallest prime
divisor of $|G|$. Then the following inequalities are true:
$None$
$\frac{d(H,K)}{|J(G,H,K)|}+\frac{|C^{\otimes}_{K}(H)|}{|H|}\left(1-\frac{1}{|J(G,H,K)|}\right)\leq
d^{\otimes}(H,K)$ $None$ $d^{\otimes}(H,K)\leq
d(H,K)-\left(1-\frac{1}{p}\right)\left(\frac{|C_{K}(H)|-|C^{\otimes}_{K}(H)|}{|H|}\right).$
On the other hand, we may correlate the relative tensor degree, the relative
commutativity degree and another notion, studied recently in [11]. In order to
proceed in this direction, we recall from [3, 5, 10] that the nonabelian
exterior product $H\wedge K$ of $H$ and $K$ is the quotient the nonabelian
tensor produc $H\otimes K$, defined by $H\wedge K=(H\otimes K)/\nabla(H\cap
K)=\langle(x\otimes y)\nabla(H\cap K)\ |\ x,y\in H\cap K\rangle=\langle
x\wedge y\ |\ x,y\in H\cap K\rangle$, where $\nabla(H\cap K)=\langle x\otimes
x|\ x\in H\cap K\rangle$. From [3, 4], we may note that
$\kappa^{\prime}_{H,K}:h\wedge k\in H\wedge
K\mapsto\kappa^{\prime}_{H,K}(h\wedge k)=[h,k]\in[H,K]$
is an epimorphism of groups such that
$\begin{CD}1@>{}>{}>M(G,H,K)@>{}>{}>H\wedge
K@>{\kappa^{\prime}_{H,K}}>{}>[H,K]@>{}>{}>1\end{CD}$
is a central extension, where $M(G,H,K)=\ker\kappa^{\prime}_{H,K}$ is the
so–called Schur multiplier of the triple $(G,H,K)$. We inform the reader that
several references on the theory of the Schur multipliers of triples can be
found in [4, 11]. In particular, $M(G,G,G)=M(G)=H_{2}(G,\mathbb{Z})$ is the
Schur multiplier of $G$, that is, the second integral homology group over $G$.
In our situation, it is possible to consider the set
$C^{\wedge}_{K}(H)=\\{k\in K\ |\ h\wedge k=1,\ \forall h\in H\\}=\bigcap_{h\in
H}C^{\wedge}_{K}(h),$
called exterior centralizer of $H$ with respect to $K$ and it is actually a
subgroup of $K$ (see [10] for details). In particular,
$C^{\wedge}_{G}(G)=Z^{\wedge}(G)=\bigcap_{x\in G}C^{\wedge}_{G}(x)$ is called
exterior center of $G$. It is easy to check that $C^{\wedge}_{G}(x)\subseteq
C_{G}(x)$ and $Z^{\wedge}(G)\subseteq Z(G)$.
Some recent papers as [11] show that it is possible to have a combinatorial
approach for measuring how far a group $G$ is from $Z^{\wedge}(G)$ and this is
interesting, because a result of Ellis [5] characterize a capable groups by
the triviality of its exterior center (i.e.: a group $G$ is capable if
$G\simeq E/Z(E)$ for a given group $E$). This aspect has motivated the notion
of relative exterior degree
$d^{\wedge}(H,K)=\frac{|\\{(h,k)\in H\times K\ |\ h\wedge
k=1\\}|}{|H||K|}=\frac{1}{|H|\ |K|}\sum_{h\in H}|C^{\wedge}_{K}(h)|$
of $H$ and $K$. When $G=H=K$, we find the exterior degree
$d^{\wedge}(G,G)=d^{\wedge}(G)$ of $G$ in [11]. It is easy to prove that
$d^{\wedge}(G)=1$ if and only if $G=Z^{\wedge}(G)$. Hence the exterior degree
represents the probability that two randomly chosen elements commute with
respect to the operator $\wedge$. Roughly speaking, this means that there are
many chances of finding capable groups for small values of exterior degree.
From [12, Theorem 2.8], we may correlate the above notions via the inequality
$d^{\otimes}(G)\leq d^{\wedge}(G)\leq d(G)$
and our second main theorem shows that something of more general holds.
###### Theorem 1.2.
Let $H,K$ be normal subgroups of a group $G$. Then
$d^{\otimes}(H,K)\leq d^{\wedge}(H,K)\leq d(H,K).$
Moreover, if $J(G,H,K)$ is trivial, then
$d^{\otimes}(H,K)=d^{\wedge}(H,K)=d(H,K)$.
## 2\. Proofs of the main results
We begin with a technical lemma, whose proof uses an argument which appears in
[11, Lemma 2.1] and [12, Lemma 2.2] in different ways.
###### Lemma 2.1.
Let $H,K$ be normal subgroups of a group $G$. Then
$d^{\otimes}(H,K)=\frac{1}{|H|}\sum^{k_{K}(H)}_{i=1}\frac{|C^{\otimes}_{K}(h_{i})|}{|C_{K}(h_{i})|}.$
In particular, if $G=HK$, then $C_{K}(h_{i})/C^{\otimes}_{K}(h_{i})$ is
isomorphic to a subgroup of $J(G,H,K)$ and
$|C_{G}(h_{i}):C^{\otimes}_{G}(h_{i})|\leq|J(G,H,K)|$ for all
$i=1,2,\ldots,k_{K}(H)$.
###### Proof.
Since $H$ is normal in $G$, we consider the $K$–conjugacy classes
$C_{1},\ldots,C_{k_{K}(H)}$ that constitute $H$. It follows that
$|H|\ |K|\ d^{\otimes}(H,K)=\sum_{h\in
H}|C^{\otimes}_{K}(h)|=\sum^{k_{K}(H)}_{i=1}\sum_{h\in
C_{i}}|C^{\otimes}_{K}(h)|$ $=\sum^{k_{K}(H)}_{i=1}|K:C_{K}(h_{i})|\
|C^{\otimes}_{K}(h_{i})|=|K|\
\sum^{k_{K}(H)}_{i=1}\frac{|C^{\otimes}_{K}(h_{i})|}{|C_{K}(h_{i})|}.$
Now assume that $G=HK$. For all $i=1,\ldots,k_{K}(H)$, the map
$\varphi:kC^{\otimes}_{K}(h_{i})\in
C_{K}(h_{i})/C^{\otimes}_{K}(h_{i})\longmapsto k\otimes h_{i}\in J(G,H,K)$
satisfies the condition
$\varphi(k_{1}k_{2}C^{\otimes}_{K}(h_{i}))=k_{1}k_{2}\otimes
h_{i}=(k_{1}\otimes h_{i})^{k_{2}}\ (k_{2}\otimes h_{i})$ $=(k_{1}\otimes
h_{i})\ (k_{2}\otimes h_{i})=\varphi(k_{1}C^{\otimes}_{K}(h_{i}))\
\varphi(k_{2}C^{\otimes}_{K}(h_{i}))$
for all $k_{1},k_{2}\in C_{K}(h_{i})$. Furthermore,
$\ker\varphi=\\{kC^{\otimes}_{K}(h_{i})\ |\ k\otimes
h_{i}=1\\}=C^{\otimes}_{K}(h_{i})$. Then $\varphi$ is a monomorphism and
$C_{K}(h_{i})/C^{\otimes}_{K}(h_{i})$ is isomorphic to a subgroup of
$J(G,H,K)$. We conclude that
$|C_{K}(h_{i}):C^{\otimes}_{K}(h_{i})|\leq|J(G,H,K)|$. ∎
Now we may prove Theorem 1.1. It is an interesting bound, which connects the
notion of relative tensor degree with that of relative commutativity degree.
###### Proof of Theorem 1.1.
We begin to prove $(a)$. From Lemma 2.1, we have
$|C^{\otimes}_{K}(H)|/|C_{K}(H)|\geq 1/|J(G,H,K)|$
and, together with the equality $d(H,K)=\frac{k_{K}(H)}{|H|}$, we deduce
$d^{\otimes}(H,K)=\frac{1}{|H|}\sum^{k_{K}(H)}_{i=1}\left|\frac{C^{\otimes}_{K}(h_{i})}{C_{K}(h_{i})}\right|\geq\frac{1}{|H|}\left(|C^{\otimes}_{K}(H)|+\frac{k_{K}(H)-|C^{\otimes}_{K}(H)|}{|J(G,H,K)|}\right)$
$=\frac{k_{K}(H)}{|H|\
|J(G,H,K)|}+\frac{|C^{\otimes}_{K}(H)|}{|G|}\left(1-\frac{1}{|J(G,H,K)|}\right)$
$=\frac{d(H,K)}{|J(G,H,K)|}+\frac{|C^{\otimes}_{K}(H)|}{|H|}\left(1-\frac{1}{|J(G,H,K)|}\right)$
Conversely, $|K:C^{\otimes}_{K}(h_{i})|\geq p$ implies that
$d^{\otimes}(H,K)=\frac{1}{|H|}\sum^{k_{K}(H)}_{i=1}\left|\frac{C^{\otimes}_{K}(h_{i})}{C_{K}(h_{i})}\right|$
$\leq\frac{|C^{\otimes}_{K}(H)|}{|H|}+\frac{1}{p}\left(\frac{|C_{K}(H)|-|C^{\otimes}_{K}(H)|}{|H|}\right)+\frac{k_{K}(H)-|C_{K}(H)|}{|H|}$
$=d(H,K)-\frac{p-1}{p}\left(\frac{|C_{K}(H)|-|C^{\otimes}_{K}(H)|}{|H|}\right).$
∎
Immediately, we note that [12, Theorem 2.3] describes a special case of
Theorem 1.1. The following consequence of Theorem 1.1 is interesting, too.
###### Corollary 2.2.
Let $G=HK$ for two normal subgroups $H$ and $K$. Then
$\frac{d(H,K)}{|J(G,H,K)|}\leq d^{\otimes}(H,K)\leq\ d(H,K).$
In particular, if $J(G,H,K)$ is trivial, then $d^{\otimes}(H,K)=d(H,K).$
The second main theorem is a result of comparison. Its proof is the following.
###### Proof of Theorem 1.2.
We have
$d^{\wedge}(H,K)=\frac{1}{|H|}\sum^{k_{K}(H)}_{i=1}\left|\frac{C^{\wedge}_{K}(h_{i})}{C_{K}(h_{i})}\right|\leq\frac{1}{|H|}\sum^{k_{K}(H)}_{i=1}\left|\frac{C_{K}(h_{i})}{C_{K}(h_{i})}\right|=d(H,K)$
and the upper bound follows.
Now $k\in C^{\wedge}_{K}(H)$ if and only if $k\wedge h=1$ for all $h\in H$ if
and only if $(k\otimes h)\nabla(H\cap K)=\nabla(H\cap K)$ if and only if
$k\otimes h\in\nabla(H\cap K)$. This condition is weaker than the condition
$k\otimes h=1$, characterizing the elements of $C^{\otimes}_{K}(H)$. Then
$C^{\otimes}_{K}(H)\subseteq C^{\wedge}_{K}(H)\subseteq C_{K}(H)$. This and
Lemma 2.1 imply the lower bound
$d^{\otimes}(H,K)=\frac{1}{|H|}\sum^{k_{K}(H)}_{i=1}\left|\frac{C^{\otimes}_{K}(h_{i})}{C_{K}(h_{i})}\right|\leq\frac{1}{|H|}\sum^{k_{K}(H)}_{i=1}\left|\frac{C^{\wedge}_{K}(h_{i})}{C_{K}(h_{i})}\right|=d^{\wedge}(H,K).$
The rest follows from Corollary 2.2. ∎
## References
* [1] A.M.A. Alghamdi and F.G. Russo, A generalization of the probability that the commutator of two group elements is equal to a given element, Bull. Iranian Math. Soc. 38 (2012), 973–986.
* [2] A. Amit and U. Vishne , Characters and solutions to equations in finite groups, J. Algebra Appl. 10 (2011), 675–686.
* [3] R. Brown, D. L. Johnson and E. F. Robertson, Some computations of non-abelian tensor products of groups, J. Algebra 111 (1987), 177–202.
* [4] R. Brown, P.J. Higgins and R. Sivera, Nonabelian algebraic topology, EMS Tracts in Mathematics, EMS Publishing, Zürich, 2011.
* [5] G. Ellis, Tensor products and $q$–crossed modules, J. London Math. Soc. 51 (1995), 243–258.
* [6] R.M. Guralnick and G.R. Robinson, On the commuting probability in finite groups, J. Algebra 300 (2006), 509–528.
* [7] W.H. Gustafson, What is the probability that two groups elements commute? Amer. Math. Monthly 80 (1973), 1031–1304.
* [8] P. Lescot, Isoclinism classes and commutativity degrees of finite groups, J. Algebra 177 (1995), 847–869.
* [9] P. Lescot, Central extensions and commutativity degree, Comm. Algebra 29 (2001), 4451–4460.
* [10] P. Niroomand and F.G. Russo, A note on the exterior centralizer, Arch. Math. (Basel) 93 (2009), 505–512.
* [11] P. Niroomand, R. Rezaei and F.G. Russo, Commuting powers and exterior degree of finite groups, J. Korean Math. Soc. 49 (2012), 855–865.
* [12] P. Niroomand and F.G. Russo, On the tensor degree of finite groups, Ars Comb., to appear. Available as preprint at: http://arxiv.org/abs/1303.1364.
* [13] D.E. Otera, F.G. Russo and C. Tanasi, Some algebraic and topological properties of nonabelian tensor product, Bull. Korean Math. Soc. 50 (2013),1069–1077.
|
arxiv-papers
| 2013-09-23T21:13:56 |
2024-09-04T02:49:51.396576
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Ahmad M.A. Alghamdi (Umm Alqura University, Mekkah, Saudi Arabia) and\n Francesco G. Russo (Universita' degli Studi di Palermo, Palermo, Italy)",
"submitter": "Francesco G. Russo",
"url": "https://arxiv.org/abs/1309.5983"
}
|
1309.5984
|
# An evolutionary approach to Function
Phillip Lord
School of Computing Science
Newcastle University
Newcastle-Upon-Tyne
United Kingdom
NE1 7RU
[email protected]
###### Abstract
Background: Understanding the distinction between function and role is vexing
and difficult. While it appears to be useful, in practice this distinction is
hard to apply, particularly within biology.
Results: I take an evolutionary approach, considering a series of examples, to
develop and generate definitions for these concepts. I test them in practice
against the Ontology for Biomedical Investigations (OBI). Finally, I give an
axiomatisation and discuss methods for applying these definitions in practice.
Conclusions: The definitions in this paper are applicable, formalizing current
practice. As such, they make a significant contribution to the use of these
concepts within biomedical ontologies.
## 1 Background
Large parts of modern biology are aimed at answering questions about function.
For example, much of the Gene Ontology deals with molecular function(Ashburner
et al., 2000). In dealing with the social aspects of science, roles are
similarly important. It is clear, therefore, that function and role are
important concepts in biomedical ontologies and are prime candidates for
inclusion in an upper ontology. A coherent, consistent and shared definition
for function and role is likely to decrease the effort required to integrate
independently-developed ontologies.
One upper ontology, the Basic Formal Ontology (BFO)(bfo, ), currently in use
by a number of groups, includes a definition of function. However, this
definition is not naturally applicable to biology; it is not clear, for
instance, that Gene Ontology molecular functions are also BFO functions, as
discussed later (Section 2.1). An alternative is available within the General
Formal Ontology (GFO), which provides an ontology of function(Burek et al.,
2006). This ontology provides a more extensive framework for describing
functions but, in itself, does not define biological function. Despite this,
there is a reasonable degree of informal agreement among biologists as to the
meaning of the word. Conversely, while formal and informal definitions for
role seems clear, many people have difficulties in applying it in practice.
In this paper, I address two key issues relating to the modeling of function
and role coherently for biomedical ontologies: firstly, how I unify the
definitions of function as they apply to artifacts and to life; and, secondly,
how do I differentiate between roles and functions? I consider illustrative
examples where the answers are reasonably clear and evolve a definition from
them. I then consider examples of the application of current definitions from
BFO in a practical biomedical use. Finally, I offer an axiomatisation in OWL
stemming from my definitions and consider how they could be applied in
practice.
## 2 Results
### 2.1 Biological Function
First, consider the current definition of function provided by BFO (see Def:
1).
###### Definition 1
Function is a realizable entity the manifestation of which is an essentially
end-directed activity of a continuant entity being a specific kind of entity
in the kind or kinds of context that it is made for
As a simple example, a hammer (_the continuant entity_) was made to hammer
nails (_the function_) in a hammering process (_the end-directed activity_).
This definition is problematic for biological systems. The problem here is
simple: most biological systems were not made or designed for any purpose.
Although, it has not been incorporated into BFO yet, there is a potential
definition for “biological function” which would become a new child of
function (Def: 2)(Arp and Smith, 2008).
###### Definition 2
A biological function is a function which inheres in an independent continuant
that is i) part of an organism and ii) exists and has the physical structure
as a result of the coordinated expression of that organism’s structural genes.
As an example, a foot (_the independent continuant_) is part of an organism,
exists, and develops in a controlled way as a result of gene expression. There
are some difficulties with this definition. Consider the following examples: a
differentiated tumour has its structure through coordinated expression of its
genes, it exists and it engages in end-directed activity; a male ant has its
structure as a result of gene expression, and engages in end-directed
activity, however, it is not part of an organism. From this, I conclude that a
differentiated tumour does, indeed, have a function (growing). A male ant,
however does not have a function. Also consider molecular function: the
physical structure of a protein is independent of the expression of an
organisms structural genes – only its presence depends on this. A protein,
therefore, does not have a function, by this definition. So, there are two key
problems: the definition does not work for entities above or below a certain
size; and most biological entities have their structure as a result of
coordinated expression. I offer the following alternate definition (Def: 3).
###### Definition 3
A biological function is a realizable entity that inheres in a continuant
which is realized in an activity, and where the homologous structure(s) of
individuals of closely related and the same species bear this same biological
function.
The definition given uses the notion of homology; evolution is key to our
understanding of biology and it is appropriate that it should be used to
define biological function. If a biological structure has a function, then
this function will have evolved along with the structure; so, other structures
with a common evolutionary descent will display the same behaviour. This
definition also mirrors closely normal biological practice; the most common
way to determine the function of an unknown structure is to look for function
of a homologous structure.
It should be noted that this definition of biological function is _not_
circular, although it has itself as part of its definition; rather it is
recursive; a chimp hand and a human hand can have the same function because of
each other. It does require that a structure must have a homolog for it to
have a function. It does not require that these homologous structures be
extant.
Applying this to our examples: the tumour now has no function because it has
no homologs (different tumours arise as independent events and share no common
ancestor). Likewise, the activity of the male ant now clearly is a function,
as many different, related organisms behave in a similar way. Finally, a
protein may have a function depending on the activity of its homologs. In
short, this definition results in the same conclusions as our biological
understanding.
### 2.2 Relating the Functions
As well as biological function, another subclass, artifactual function has
been suggested(Arp and Smith, 2008). Next, I consider the relationship between
function and these biological and artifactual subclasses. Taking an
illustrative example, consider the sole of my foot and the sole of my shoe.
They appear to operate to: provide a frictional surface to enable motion;
provide padding to reduce shock to everything above; be tough enough to resist
abrasion. They would appear to have the same function; indeed, like many
artifacts, we would guess that the shoe owes much of its design to mimicry of
biology. It would seem, therefore, sensible that instances of Shock Resistance
could be either a biological or artifactual function. The alternative would be
to duplicate many functions under both subclasses (“biological shock
resistance”, and “artifactual shock resistance”) with very similar
definitions. Some functions such as reproduction have to be a biological
function, while others liquifying iron can never be. In short, whether an
instance of function is biological or artifactual should be determined from
the nature of the entity in which it inheres, rather than the process by which
it is realized. I therefore offer a simple definition of function which
reflects this (Def: 4).
###### Definition 4
A function is a realizable entity which is a biological function or an
artifactual function.
For the purposes of this paper, I note that the definition given for function
earlier (Def: 1), can serve as a reasonable definition for artifactual
function.
It is also interesting that this definition covers some unusual but
nonpathological examples. Take a bacterium whose colonies change colour
depending on the presence of a toxin and which was produced using synthetic
biology techniques. The components have all evolved, but the organisation has
not. Is the detection of the toxin then a biological or artifactual function?
This is clearly a difficult, if uninspiring, question but given Def: 4, one
that can avoid by simply describing it as a function; perhaps more
intuitively, it can be described as both a biological and an artifactual
function, suggesting strongly that these two classes should not be disjoint.
### 2.3 Roles
Next, having considered the definition of function and its applicability to
biology, I consider the issue of roles. The current definition (Def: 5) is
complex; put more simply, it suggests that the entity having that role can be
involved in an activity but that it was not necessarily intended for, nor
necessarily has the structure for this.
###### Definition 5
A realizable entity the manifestation of which brings about some result or end
that is not essential to a continuant in virtue of the kind of thing that it
is but that can be served or participated in by that kind of continuant in
some kinds of natural, social or institutional contexts.
Consider the relationship between role and function. Again, I shall use a
simple biological example, in this case of a man walking on his hands. By our
earlier definition (Def: 3), “to walk on hands” is _not_ a function. While the
homologous structure is, indeed, used for walking on in all closely related
species, most humans do not walk on their hands. It would, therefore appear to
be a role. In this context the hand has a role of Shock Resistance. This
realizable entity also appears in the hands of many other primates; in this
case, however, it would appear to be a function of the primate hand, as it is
a function of their feet.
I am left with a similar conclusion as previously. Just as Shock Resistance
maybe either a biological or artifactual function, I must also conclude that
the individuals of the same class can be a role, depending on the nature of
the relationship between the independent continuant and realizable entity.
### 2.4 OBI as a case study
So far in this paper, I have considered a number of illustrative examples and
used these to draw conclusions about definitions for functions and roles. This
methodology is appropriate, but has the limitation that the choice of other
examples may have led us to different conclusions. In this section, therefore,
I will consider the use that OBI (Ontology for Biomedical Investigations)(The
OBI consortium, 2009) has made of function and role (Analysis was performed on
OBI rc-1, (release 2009-11-06)). I choose to use OBI as it was built after BFO
and with knowledge of it; many of the ontologies available from OBO were
started without its use or knowledge.
Function Bearer | OBI Function
---|---
Human | Perturb, Measure, Separation (3), Sort
Computer | Information Processor (3), Consume Data
Highly Generic | Freeze, Heat, Environment Control, Mechanical, Record, Contain, Transfer, Cool, Connection, Synthesizing, Excitation (2), Ionization, Energy Supply (1)
Distant Galaxy | Magnify
Out of Scope | Molecular Function (3)
Table 1: OBI Functions considered as Roles. I provide suggestions for entities
that could engage in the same end-directed activity, but without being
designed for the purpose. Function has been omitted from the OBI term names.
Numbers indicate child terms which have been omitted for brevity. I do not
consider Information Processor to be the function of a computer, as the
definition is more specific than the term suggests.
Considering first the functions of which OBI has 38. Can these functions be
fulfilled by an entity which was not designed for the purpose? As shown in
Table 1 most of them can, often by considering a highly generic device (like a
computer) or organism (like a human). Some are highly generic in themselves
and can be performed by many things (heat for example). I find no cases where
a function could not be fulfilled by an entity which was not designed for the
purpose. I consider “Molecular Function” to be out-of-scope for this section
as it is not clear whether it fulfils the BFO definition of function.
Next, I consider OBI roles. There are many more roles than functions – around
90 in fact. Due to the size, here I consider whole branches of the role
hierarchy.
There are a set of roles relating to reagents and their states. While label
role (defined as a reagent role realized in a detection of label assay) seems
sensible, it would appear that S35 CTP, used to label DNA, is manufactured
specifically for this purpose. It certainly does not occur in nature. It would
appear to fulfil the requirements for a function.
Next, let us consider reference substance roles (a role which can support the
observation of similarities, differences, relative magnitude or change). In
many cases, biological assays use a reference which is not manufactured.
However, consider $\lambda$-HindIII fragments or a calibration standard. These
would all appear to function as an reference and have been produced
specifically for this purpose.
Finally, there are a number of roles for molecules or organisms: antigen role,
pathogen role and primer role. In an age where we can engineer the production
of DNA, protein and organisms, it is not clear that these can only ever be
roles and not functions.
As a result of this analysis, I suggest a modification to the current
definition of role (Def: 5). Both roles and functions can become apparent
(_realized_) in natural, social or institutional contexts. That such a context
exists does not provide a clear differentation between role and function; the
critical distinction relates to whether the entity in question was designed to
be or has homologs that are engaged in a given process. I suggest, therefore,
this alternate and simpler definition for role (Def: 6).
###### Definition 6
A realizable entity the manifestation of which brings about some result or end
that is not essential to a continuant because of its kind.
In short, from this case study, I conclude that the role/function distinction
is not clear. While OBI has a specific intent in mind with its application of
the distinction (broadly and not exhaustively, social or experimental roles,
device or instrument functions), this distinction is not the distinction made
in the current definitions of role and function in BFO; further given that
most functions could also appear to be roles, and many roles appear also to be
functions, I suggest that the distinction made in the current definitions is
not useful in the context of OBI.
The earlier theoretical analysis seems to be confirmed in practice within OBI.
This suggests that the limitations in the definitions drawn from the earlier
illustrative examples are general and not simply as a result of the specific
examples chosen.
### 2.5 An axiomatisation for function and role
I have produced an axiomatisation in OWL of functions and roles as defined in
this paper, available as described in the abstract; due to space
considerations I report here key differences between this and the BFO
axiomatisation.
* •
There is an explicit relationship between RealizableEntity and Process.
Subclasses use a more specific relationship. So, a ToAbsorbShock function may
only be realized by a AbsorbtionOfShock process, if it is realized at all.
* •
Function and Role are defined classes. Stating that an instance of
ToAbsorbShock is_function_of instance of FootSole implies, therefore, that the
former is a function.
* •
Most leaves of RealizableEntity are direct children of RealizableEntity, with
a few exceptions (ToReproduce is a child of BiologicalFunction).
The definitions could be extended further; for simplicity, I have not added
classes to differentiate between organisms and artifacts. These could be added
to automate the population of BiologicalFunction and ArtifactualFunction.
At the current time, it remains an open question whether Role, Function and
its children should be disjoint. The key example of the function of a
synthetic biological organism suggests that Biological and ArtifactualFunction
function should not be disjoint (as it appears to be both), but I have no
example which suggests whether Role and Function should be disjoint. In
axiomatisating the examples given, these disjoint statements make no practical
difference.
Many ontologies are built using the OBO format; while this has a slightly
weaker semantics than OWL it is possible to represent much of OWL in OBO
format(Golbreich et al., 2007). The axiomatisation presented here can be
represented using OBO union and intersection of to describe classes, which is
usually translated as a definition. The universal link between
RealizableEntity and Process has no natural equivalent. However, as this link
has its own specific relationship which is restricted to this use, problems
caused by the lack of an inexact semantic equivalent are likely to be
relatively minor.
The axiomatisation presented here is related to that produced by others;
Dumontier(Dumontier, 2008) focuses more on roles, while Burek et al.(Burek et
al., 2006) provides for a more complex representation, covering issues such as
preconditions.
### 2.6 Applying the Definitions in Practice
Finally, I consider whether these definitions are _applicable_ ; for a given
set of entities how do we decide whether we have a function (of either
subclass) or a role.
The definition of an artifactual function easily allows its application:
first, we determine whether the entity in question is an organism or part of
one (which it should not be); second, we could ask whoever produced the entity
what it was designed for. Of course, the second may not always be possible, in
which case, we can guess from its design what its purpose is. In most cases,
these questions will provide a clear answer.
For biological function, the situation is less clear. Whether an entity is an
organism or part of one is, in practice, likely to be straightforward for
extant entities; otherwise, we can apply palentological techniques. Likewise,
identification of closely related species and homologous structures is well
known as it forms the basis of taxonomy. While developing an exact definition
for “closely related” is outside the scope of this paper, it is possible.
The definition that I introduce for Function in this paper (Defn 4) is
conjunctive; it is either biological or artifactual. Here I have given little
evidence that these are the only kind of function. Fundamentally, these two
arise from very different mechanisms. There could be other appropriate
subclasses of function; the most obvious possibility would be Chemical
Function. However, artifacts are designed by humans who understand, mimic and
improve on biology by building tools. It is this mimicry that we wish to
reflect with a common definition joining biological and artifactual function;
this is not true for Chemical Function.
To determine whether something is a role, it is possible to make a
determination on the basis of whether the context is optional(Arp and Smith,
2008). However, this optionality is a difficult criterion; firstly, all
RealizableEntity’s are optional in the sense that they might never be realized
and, secondly, the optionality can depend on how specifically we define the
bearer. A hammer is not designed to hammer nails, as claimed earlier, it is
designed to hit things; a nail hammer is designed to hit nails, a toffee
hammer to hit toffee, a warhammer to hit anyone who irritates you. In
practice, a role can be considered to be a negative definition; if there is a
continuant and an end-directed activity that the continuant can be involved
in, and this involvement is known not to fulfil the definition of either
function, then we have a role.
In this paper, I have considered OBI and found that the distinction between
role and function is hard to apply; this is not true for all ontologies. For
example, consider the Gene Ontology. In many cases, the homology will be
considered as a standard part of the operating procedure(goe, ) in determining
the function of a gene product; regardless, the evidence codes would allow us
to make the distinction. We can conclude, therefore, that when the Gene
Ontology is used to annotate a protein, this describes a biological function
rather than a role.
### 2.7 Life is hard
So far, I have considered a set of examples and how the definitions might be
applied, including examples from OBI which have not been preselected. However,
categorising life is hard; here, I consider some examples which present
difficulties for the definitions I have given and the implications of these
examples.
In the first example, which I term a _drop out_ species, consider a human
walking on their hands. Earlier (Section 2.3), I have suggested that this
should be considered a role. Most of the primates do, however, walk on their
hands. However, given that the homologous structure of closely related species
use the structure for the same purpose, the definition of function (Def: 3)
would appear to apply. It is for this reason that the definition specifies
that (most) individuals of the same species must also demonstrate this
behaviour (This definition differs slightly from that given at Bio-Ontologies
2009). In short, in the absence of most individuals in a species using a
structure in a specific process, we should not use consider this structure to
have a function.
The second example, I term a _drop in_ species. Again, using a human example,
I use my larynx for vocalisation and talking. Most primates, likewise,
vocalise with their larynx; therefore, according to the given definition, this
is a function of the larynx. However, speech is considered unique to humans,
and therefore, their larynx; given that homologous structures in closely
related organisms do not bear this realizable entity, which is part of the
definition for biological function (Def: 3), I am forced to conclude that this
is a role of the larynx and not a function.
In short, while sharing a realizable entity within a species is NOT sufficient
to allow the conclusion that this entity IS a function, NOT sharing a
realizable entity within a species is sufficient to conclude that this entity
is NOT a function.
One solution to this difficulty is to state that where most individuals in a
single species use a structure within a given process, this alone is
sufficient to conclude that the structure has a function. Simply, most humans
talk with their larynx, therefore this would be a function. I counter this,
however, with the example that most humans use their fingers to operate their
mobile phones, so we would be forced to conclude that this would also be a
function. As this seems opposed to normal biological intuition and usage, I
conclude, the presence of most individuals in a species using a structure in a
specific process, is _not_ sufficient to conclude that this structure has a
function.
It is also possible that this difficulty could be resolved with greater
knowledge or changes in biology. Def: 3 does not require species be extant; if
a close, but extinct, relative of humans were shown to speak with their
larynx, or if humans speciated while maintaining their speech, again, I would
conclude that this represented a function.
While human speciation seems unlikely, it is much more relevant to other taxa.
Bacteria, in particular, evolve rapidly. There are many genes and proteins in
bacteria which are unique to a species, family or lineage(Siew et al., 2004).
In this case, the requirement for closely-related species seems to rule out
the presence of a function. Again, this seems opposed to normal biological
intuition and usage. I would counter this with two arguments. First, unlike
primates, our knowledge of the extant bacterial species is very limited. The
lack of knowledge of another speaking primate species is good evidence that no
such species exists; the lack of knowledge of a close relative for a given
bacterial species is not. Second, any definition which relies on a notion of a
species is only as good as the definition of species; for bacteria, there is
considerable debate about the utility of a species classification(Rosselló-
Mora, 2003); my definition of function will need to evolve along with our
understanding of bacterial ecology and gene flow; it may be necessary, as has
been suggested with definitions for bacterial species(Rosselló-Mora, 2003), to
have different definitions of BiologicalFunction tailored to different parts
of the taxonomy.
The evolution of a definition of role and function for proteins is difficult.
At the level of the protein, I side with Dumontier(Dumontier, 2008), who
suggests that the role/function distinction may be redundant; broadly,
proteins can do anything their structure allows, and only do things their
structure allows. The definitions given in this paper have a consistent
interpretation at the level of the protein; this avoids the necessity of
deciding at which level of granularity to stop making the role/function
distinction. We can make the distinction at all levels if we choose, but we
are not forced to do so, at those levels of granularity where it is not
useful.
It is important to note that arguing against a role/function distinction for
proteins is not to dismiss the experience of biologists in the analysis of
function assignment for genes. In this sense, the word “function” is being
used to describe an association between a protein and a process that a protein
molecule may be involved in; in short, the word “function”, in this case, can
be considered to be a synonym for “realizable entity”.
## 3 Conclusions
Here, I have taken an evolutionary approach to function and role by
considering examples and using this to derive definitions which are as
consistent as possible with current use within biomedical sciences. These
definitions have been encoded in an axiomatisation which should enable the use
of these definitions in a machine-interpretable way.
The applicability of these definitions is a key advantage; the current
distinction being made between function and role is a hard one to understand
and apply. My definition distinguishes between the two based on the nature of
the relationship to the independent continuant in which they inhere. I suggest
that it is very hard to make the distinction at the class level; my study of
OBI shows that very few of the functions and roles clearly fall into one
category or another. For an individual continuant bearing a realizable entity,
this distinction appears to be much more straightforward.
I also provide a definition of biological function, something that is
currently lacking in BFO. I have paid close attention to current biological
usage; the definition is close to the process used to determine function.
Moreover, it is highly applicable; all parts of the definition are measurable.
The desire for an applicable and measurable definition is also the reason that
I have avoided a definition based on the outcome of selective pressure; this
is hard to test in most circumstances, requiring expensive evolutionary
studies, and impossible for extinct species. Selective pressure can also be
transient. Consider industrial melanism(Majerus, 1998); should melanic
coloring be considered to gain its function during periods of pollution and
lose it in post-industrial periods? By way of analogy, should a spanner
measured in inches be considered to lose its function following metrication?
Serendipitously, it also avoids difficult questions about artificial
selection; we can state clearly that cows do not have a function of producing
beef, though this is the outcome of selection.
Importantly, my definition of biological function works across multiple levels
of granularity: from organisms and organism parts through to genes and
molecules; this is not true of previous definitions(Arp and Smith, 2008),
which cover only anatomy. It is, however, not clear how useful the
role/function distinction is for proteins and genes, as discussed earlier
(Section 2.7); It is for this reason that I have used homology rather than
orthology as the basis for the definition, as the latter is limited to the
genetic scale, where the distinction is least useful.
Finally, my definitions also do not allow distinctions that may often be made
between different types of function. For example, most biologists would
consider motion the most important function of muscle, while heat production a
byproduct; or, for a more pathological example after Hoehndorf et al. (2009),
most biologists would consider blood circulation to be a function of the
heart, but “making loud thumping noises” not to be. This is a concern which
could be best addressed by incorporating a degree of social ascription into
the categorisation of realizable entities within biology; although it is
outside the scope of this paper, this would provide a valuable and useful
addition to the current ontological practice.
In summary, I believe that the definitions and axiomatisation given in this
paper make a significant contribution to the use of role and function in
biomedical ontologies. They should enable a consistent use of these classes,
because they consider current usage of the terms and the applicability of
these definitions. I seek not to change current use but to formalize it.
## 4 Competing Interests
The author has no competing interests.
## 5 Acknowledgements
Thanks to Frank Gibson, Allyson Lister, James Malone, Helen Parkinson, Matt
Pocock and Robert Stevens for many useful discussions on the contents of this
paper. Thanks to the audience at Bio-Ontologies 2009 for the example of
rapidly-evolving bacteria.
To the memory of Mike Majerus who showed me how fundamentally strange
evolution can be.
## References
* (1) Basic Formal Ontology. URL http://www.ifomis.org/bfo.
* (2) Guide to GO Evidence Codes. URL http://www.geneontology.org/GO.evidence.shtml.
* Arp and Smith (2008) R. Arp and B. Smith. Function, role and disposition in basic formal ontology. In _The BIo-Ontologies Workshop (at ISMB 2008)_ , 2008.
* Ashburner et al. (2000) M. Ashburner, C. Ball, J. Blake, D. Botstein, H. Butler, J. Cherry, A. Davis, K. Dolinski, S. Dwight, J. Eppig, et al. Gene ontology: tool for the unification of biology. The Gene Ontology Consortium. _Nat Genet_ , 25(1):25–9, 2000.
* Burek et al. (2006) P. Burek, R. Hoehndorf, F. Loebe, J. Visagie, H. Herre, and J. Kelso. A top-level ontology of functions and its application in the open biomedical ontologies. _Bioinformatics_ , 22(14):e66–e73, Jul 2006. doi: 10.1093/bioinformatics/btl266. URL http://dx.doi.org/10.1093/bioinformatics/btl266.
* Dumontier (2008) M. Dumontier. Situational modeling: Defining molecular roles in biochemical pathways and reactions. In _OWLED 2008_ , 2008.
* Golbreich et al. (2007) C. Golbreich, M. Horridge, I. Horrocks, B. Motik, and R. Shearer. OBO and OWL: Leveraging semantic web technologies for the life sciences. _Lecture Notes in Computer Science_ , 4825:169, 2007.
* Hoehndorf et al. (2009) R. Hoehndorf, J. Kelso, and H. Herre. Contributions to the formal ontology of functions and dispositions: An application of non-monotonic reasoning. In _Bio-Ontologies 2009: Knowledge in Biology_ , 2009.
* Majerus (1998) M. Majerus. _Melanism: Evolution in Action_. Oxford University Press, 1998.
* Rosselló-Mora (2003) R. Rosselló-Mora. Opinion: the species problem, can we achieve a universal concept? _Syst Appl Microbiol_ , 26(3):323–326, Sep 2003\.
* Siew et al. (2004) N. Siew, Y. Azaria, and D. Fischer. The orfanage: an orfan database. _Nucleic Acids Res_ , 32(Database issue):D281–D283, Jan 2004. doi: 10.1093/nar/gkh116. URL http://dx.doi.org/10.1093/nar/gkh116.
* The OBI consortium (2009) The OBI consortium. Modeling biomedical experimental processes with obi. In _Bio-Ontologies 2009: Knowledge in Biology_ , 2009.
|
arxiv-papers
| 2013-09-23T21:15:10 |
2024-09-04T02:49:51.402166
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Phillip Lord",
"submitter": "Phillip Lord Dr",
"url": "https://arxiv.org/abs/1309.5984"
}
|
1309.6089
|
# Hysteretic phenomena in a 2DEG in quantum Hall effect regime studied in a
transport experiment
M. V. Budantsev [email protected] A. V. Rzhanov Institute of Semiconductor
Physics of SB RAS, Novosibirsk, Russia Novosibirsk State University, Russia
D. A. Pokhabov A. G. Pogosov E. Yu. Zhdanov A. V. Rzhanov Institute of
Semiconductor Physics of SB RAS, Novosibirsk, Russia Novosibirsk State
University, Russia A. K. Bakarov A. I. Toropov A. V. Rzhanov Institute of
Semiconductor Physics of SB RAS, Novosibirsk, Russia
###### Abstract
We investigated experimentally non-equilibrium state of a two-dimensional
electron gas (2DEG) in the quantum Hall effect (QHE) regime, studying the
hysteresis of magnetoresistance of a 2DEG with a constriction. The large
amplitude of the hysteresis enabled us to make the consistent phenomenological
description of the hysteresis. We studied the dependence on the magnetic field
sweep prehistory (minor loop measurements), recovered the anhysteretic curve,
and studied the time dependence of the magnetoresistance. We showed that the
hysteresis of magnetoresistance of a 2DEG in the QHE regime has significant
phenomenological similarities with the hysteresis of magnetization of
ferromagnetic materials, showing multistability, jumps of relaxation, and
having the anhysteretic curve. Nevertheless, we revealed the crucial
difference, manifested itself in an unusual inverted (anti-coercive) behavior
of the magnetoresistance hysteresis. The time relaxation of the hysteresis has
fast and slow regimes, similar to that of non-equilibrium magnetization of a
2DEG in QHE regime pointing to their common origin. We studied the dependence
of the hysteresis loop area on the lithographic width of the constriction and
found the threshold value of width $\sim$1.35 $\mu$m beyond which the
hysteresis is not observed. This points to the edge nature of the non-
equilibrium currents (NECs) and allows us to determine the width of the NECs
area ($\sim$0.5 $\mu$m). We suggest the qualitative picture of the observed
hysteresis, based on non-equilibrium redistribution of the electrons among the
Landau level states and assuming huge imbalance between the population of bulk
and edge electronic states.
###### pacs:
71.10.Ca, 71.70.Di, 73.43.Qt
## I Introduction
A two-dimensional electron gas (2DEG) in the quantum Hall effect (QHE) regime
demonstrates a number of intriguing phenomena that are not fully explained up
to date. Among them are non-equilibrium phenomena manifested in hysteresis of
sheet electron density, Pudalov1984 ; Dolgopolov1992 ; Dolgopolov1993
magnetization (see for review Ref.[Usher2009, ]) and electrochemical potential
Huels2004 ; Klaffs2004 as a function of magnetic field and gate voltage (for
gated structures). The hysteresis was observed in an ordinary single-layer
2DEG in nonmagnetic materials such as Si MOSFETs and AlGaAs/GaAs
heterostructures at integer and fractional filling factors corresponding to a
longitudinal resistance ($R_{\mathrm{L}}$) vanishing.
It was found that the hysteretic variations of magnetization Usher2009 and
electrochemical potential Huels2004 ; Klaffs2004 exceed equilibrium values by
a factor of 10-60. In contrast the conventional electron transport
measurements in QHE regime with vanishing $R_{\mathrm{L}}$ are not sensitive
to the hysteresis, though, the question of a probable indirect influence of
the non-equilibrium on the transport properties of 2DEG still remains open.
Recently in a number of papers it was demonstrated that creation of a
constriction in a macroscopic bath of 2DEG allows to reveal the hysteretic
behavior in magnetoresistance Budantsev2007 ; Budantsev2009 ; Budantsev22009 ;
Budantsev2011; Pioro-Ladriere2006 ; Smith2011 .
It was shown that at low temperatures (lower than 100 mK) the hysteresis
relaxes over a day Pioro-Ladriere2006 which is far beyond the time usually
spent on a conventional measurement (a magnetic field sweep) in the QHE
regime. In terms of usual magnetic hysteresis description Bertotti1998 this
type of hysteresis is considered as time-independent.
The most common explanation of the hysteretic phenomenon is based on non-
equilibrium currents (NECs) slowly relaxing providing $R_{\mathrm{L}}$
vanishing. But there are no unambiguous models in the literature, adequately
describing this phenomenon. The detailed consideration of NECs induction is
proposed in the Ref. [Shikin2002, ], but it concerns only the case of weak
non-equilibrium. Ruhe et al. Ruhe2009 studied the case of time-dependent
hysteresis. The most developed model up to now Matthews2004 is based on the
assumption of redistribution of “frozen” charges over the whole bulk area of
the 2DEG Dyakonov1991 . These charges give rise to the radial (for the disk
geometry) electric field that in turn gives rise to azimuthal drift
magnetization current. The magnetization maximum is determined by the QHE
breakdown current, calculated in the frames of quasi-elastic inter-Landau
level scattering (QUILLS) approach.Eaves1986 The mentioned model correctly
describes experimentally observed linear temperature suppression of the
hysteresis amplitude observed both in magnetization Matthews2004 and
constriction conductance.Budantsev2007 ; Budantsev2009 But the magnetization
maximum calculated in the model is less than experimentally obtained value by
a factor of 4. Moreover, the model predicts bulk NECs distribution, which, to
our opinion, contradicts with experimental data of Klaffs et. al. Klaffs2004
In the present work we study the hysteresis of the magnetoresistance of the
conventional 2DEG with constriction, placed in it. It was found that at
certain critical width of the constriction the hysteresis of the
magnetoresistance vanishes. This observation strongly suggests that NECs are
localized near the edge of the sample.
A phenomenological comparison of the 2DEG magnetoresistance hysteresis in the
QHE regime with the hysteresis of magnetization of ordinary ferromagnetic
materials is performed. Particularly, we discuss the dependence on the
magnetic field sweep prehistory for both cases, studying minor loops,
anhysteretic curves and time relaxations. Mostly we observe a similar
behavior. However significant difference has been found manifested itself in
unusual advancing (anti-coercive) response (2DEG magnetoresistance) to
external parameter (magnetic field) change, which is just opposite to the
behavior of ferromagnetic materials.
We compare our results with the hysteresis of magnetoresistance in systems
with pseudo-spin degree of freedom — so-called “quantum Hall ferromagnets”
Piazza1999 ; DePoortere2000 ; Jungwirth2001 . In some cases (see e.g. Ref.
[Piazza1999, ]) analogical anti-coercive behavior is found. This
phenomenological similarity suggests that there are similar mechanisms
underlying the origin of the hysteresis both in ordinary 2DEG and in QH
ferromagnets.
In addition to the phenomenological description, we provide a qualitative
picture of a spatial distribution of NECs and electrochemical potential in the
sample.
The paper is organized in the following way. In Sec. II we describe the
experimental detail. In Sec. III we present the experimental data of following
measurements: the magnetoresistance of the 2DEG with the constriction in the
QHE regime, minor loop, anhysteretic curve, time relaxation dependence of the
magnetoresistance, and the dependence of the magnetoresistance hysteresis loop
area on the lithographic width of the constriction and discuss them. In Sec.
IV we develop a detailed qualitative picture of spatial distribution of NECs
in a sample, explaining the hysteresis of the magnetoresistance of the
constriction as well as the non-equilibrium magnetization of a 2DEG in the QHE
regime. We discuss inverted (anti-coercive) behavior of the hysteresis and
compare the hysteresis of the magnetoresistance with the hysteresis of the
ferromagnetic materials magnetization. We conclude in Sec. V.
## II Experimental details
Experimental samples were fabricated from GaAs/AlGaAs heterojunctions with
2DEG of two types (type I and type II) grown in different cycles of molecular
beam epitaxy. The electron mobility $\mu$ and spacer layer thickness $d$ are
$\mu$=0.6$\div$0.8$\times$106 cm2/V$\cdot$s and $d=400$ Å in the type I sample
and $\mu$=0.8$\div$1.0$\times$106 cm2/V$\cdot$s and $d=300$ Å in the type II
sample. The electron densities of the macroscopic 2DEGs of both types are
ns=3$\times$1011 cm-2 at the temperature of 4.2 K. Hall bars of the dimensions
$W$$\times$$L$=50$\times$100 $\mu$m2 have been created on the surface of the
heterostructures by means of photolithography. In the middle part of the Hall
bars constrictions of the effective width 0.6$\div$0.8 $\mu$m and the length 3
$\mu$m have been created by means of electron beam lithography (see the inset
of Fig. 1). The samples of the both type allow to measure longitudinal
magnetoresistance $R_{\mathrm{L}}=U_{35}/I_{12}$ and Hall magnetoresistance of
macroscopic 2DEG $R_{\mathrm{H}}=U_{36}/I_{12}$. The measurements are carried
out in the linear response regime on the alternating current of the magnitude
1$\div$10 nA and the frequency 7 Hz at the temperature of 60 mK. The magnetic
field covering a range of 0$\div$15 T was oriented perpendicular to the 2DEG.
Figure 1: The magnetoresistance of the samples of the type I. At the centers
of the plateaux of macroscopic 2DEG Hall resistance
$R_{\mathrm{H}}~{}=~{}U_{36}/I_{12}$ at odd and even integer filling factors a
giant hysteresis of longitudinal magnetoresistance
$R_{\mathrm{L}}~{}=~{}U_{35}/I_{12}$ is observed. The geometry of the sample
is shown in the inset.
## III Experimental results
First, we have studied the type I samples. A giant hysteresis of longitudinal
magnetoresistance $R_{\mathrm{L}}$ of the constriction with respect to
magnetic field sweep direction is observed in magnetic fields corresponding to
centers of the plateau of macroscopic 2DEG Hall resistance $R_{\mathrm{H}}$ at
both odd and even integer filling factors (Fig. 1). The effect is independent
on a direction of magnetic field. The most pronounced hysteresis loop is
observed near the filling factor $\nu=1$ in the range of filling factors
between $\nu_{\mathrm{h1}}=0.95$ and $\nu_{\mathrm{h2}}=1.07$ on the
background of macroscopic 2DEG Hall plateau observed in the range of filling
factors between $\nu_{1}=0.85$ and $\nu_{2}=1.22$ (Fig. 1). The hysteresis has
the amplitude 10 k$\Omega$ amounting to 100$\%$ of measuring signal and the
width width of the hysteresis is 1.2 T at filling factor $\nu=1$. Drastic
step-like rise of the magnetoresistance from zero up to 10 $k\Omega$ at up-
sweep of the magnetic field and sharp drop down to zero at down-sweep of the
magnetic field allows us to consider this behavior as the magneto-induced
breakdown of QHE.
The large amplitude of the hysteresis made it possible to carry out a detailed
study. It has been found that the response of the system (magnetoresistance)
outpaces the external excitation (magnetic field change) rather than retards
from it, unlike the case of ferromagnetic magnetization. In mathematical terms
the hysteresis loop is negatively oriented curve, while in ferromagnetic
materials the hysteresis loop is positively oriented. One can see that the
magnetoresistance dramatically changes at the entrance to the area of
hysteresis, that is $|dR/dB|$ tends to a huge value, and it remains almost
unchanged when leaving the area. In the case of retarded behavior the
magnetoresistance would practically not be changed at the entrance to the area
of hysteresis, that is $|dR/dB|$ would be close to zero. In this sense the
observed hysteresis is abnormally inverted having a negative coercivity in
contrast to the ferromagnetic magnetization.
Figure 2: The hysteresis of the magnetoresistance of the II type sample at
the filling factor $\nu=1$. (a) Major and minor loops. (b) Anhysteretic curve.
Arrows indicate the magnetic field sweep direction. The sweep-rate is 0.05
T/min. The geometry of the sample and relaxation are shown in the insets. One
can see the fast relaxation regime with $\tau$=2.4 s. The relaxation is
occurred by jumps simultaneously in both structures.
The magnetoresistance of the type II sample also demonstrates the hysteresis
of magnetoresistance (Fig. 2) in the range of filling factors between
$\nu_{\mathrm{h1}}=0.96$ and $\nu_{\mathrm{h2}}=1.06$ on the background of
macroscopic 2DEG Hall plateau observed in the range of filling factors between
$\nu_{1}=0.88$ and $\nu_{2}=1.16$. One can see that the hysteresis of
magnetoresistance of the type II samples is observed at the same range of
filling factors as in the type I samples but it does not demonstrate a giant
amplitude and drastic jumps. This distinction is probably resulted from the
differences in heterostructures parameters and needs to be discussed
separately. Nevertheless, we succeeded to reveal inverted (anti-coercive)
behavior of the magnetoresistance hysteresis of the type II sample studying
the form of minor hysteretic loops. One can see that a change of the magnetic
field sweep direction inside the major hysteretic loop leads to dramatic
change of the magnetoresistance, i.e. the derivative $|dR/dB|$ tends to a huge
value (Fig. 2(a)), while the retarding behavior would correspond to zero
$|dR/dB|$ (as in the case of minor loops of ferromagnetic materials, where
$|dM/dH|=0$). Moreover the analysis of the minor loops allows us to conclude
that there is a large number of sample states inside the major hysteresis
loop, and therefore the sample state is multistable.
It is generally accepted that long-lived NECs are induced on the background of
Hall plateaus where $\sigma_{xx}$ is vanishing. Our measurements shows that
the range of existence of the hysteresis $\Delta\nu_{\mathrm{h}}$ is
approximately 3 times less than the Hall plateau width $\Delta\nu$:
$\Delta\nu_{\mathrm{h}}$/$\Delta\nu$$\approx$1/3. That is the condition
$\sigma_{xx}\approx 0$ is necessary but not sufficient to observe the
hysteresis.
To perform consistent phenomenological comparison of the 2DEG
magnetoresistance hysteresis in QHE regime with the hysteresis of
magnetization of ferromagnetic materials we have measured the anhysteretic
curves (Fig. 2(b)). Each experimental point $R_{i}(B_{i})$ has been obtained
by cycling the magnetic field around $B_{i}$ with a decreasing amplitude. The
form of the anhysteretic curve and its position inside the loop are
phenomenologically similar to that of ferromagnetic materials.
We have studied the time relaxation of the hysteresis, measuring the time
dependence of the magnetoresistances of two conducting channels situated in
the same Hall bar and separated by macroscopic 2DEG reservoir of the length of
250 $\mu$m and the width of 50 $\mu$m (see the top inset of Fig. 2(a)). The
dependencies are shown in the insets of Fig. 2(a) and (b). The curves have
been obtained after stopping the up-sweep of the magnetic field at the value
12.25 T, corresponding to the filling factor $\nu=1$ where the hysteresis has
the maximum amplitude.
It has been found that time relaxation has two phases: the fast initial one
followed by slow phase. After stopping the magnetic field sweep we have
observed the fast exponential relaxation by relatively small value (less than
25$\%$) with relaxation time 2.4 s (see the bottom inset of Fig. 2(a)). Such
relaxation time is too big to be explained by spin-orbital coupling mechanism
of relaxation Muller1992 ; Khaetskii1992 and resembles the nuclear relaxation
time in GaAs which is about 30 s.Dixon1997 ; Devyatov2004 It should be noted
that the non-equilibrium magnetization relaxes in the similar way.Kershaw2007
It also has regime of initial exponential decay followed by a much slower
power-law decay, however the relaxation time in the first regime is longer
($\sim$20 s) than obtained in our work. In Ref. [Matthews2004, ] it has been
suggested that magnetic field change results in QHE breakdown accompanied by
electron transitions between adjacent Landau levels with opposite spin. The
observed fast relaxation can be explained by such inter-Landau levels
scattering and brings the system to a local (not global) minimum of energy
after stopping the magnetic field sweep. According to Ref. [Matthews2004, ]
such scattering results from QUILLS process Eaves1986 and takes place in a
part of the sample where the maximum electric field is reached. In Refs.
[Matthews2004, , Dyakonov1991, ] it is suggested that maximum electric field
is reached in the bulk of a sample. However our results are sensitive to the
sample edge (see Section IV) since the magnetoresistance is defined by the
edge states filling. This allows us to conclude that QHE breakdown takes place
at the sample edge.
Further relaxation is much slower. During the next 40 min the resistances of
both structures relax by a small value ($\sim$20$\%$). Moreover the most
notable changes are occurred by sudden jumps simultaneously in both remote
structures (see the inset of Fig. 2(b)). Taking into account that the
structures are separated by macroscopic 2DEG reservoir of the length of 250
$\mu$m and the width of 50 $\mu$m it can be concluded that these jumps are
caused by relaxation process in the macroscopic 2DEG reservoir. This
observation once again confirms the conclusion that the narrow conductive
channel is a tool for the study of the non-equilibrium phenomena taking place
in a macroscopic 2DEG.Budantsev2007 The observed stepwise relaxation is
phenologically similar to Barkhausen jumps observed in the ferromagnetic
materials and originating from spin domain structure
transformations.Bertotti1998
A formation of a spatial spin polarization in a single quantum wire has been
considered in Ref. [Ihnatsenka2007, ] and it has been shown that this would
result in hysteresis in electron transport measurements. However, all the
calculated characteristics demonstrate the conventional coercive behavior.
Consequently, the spin mechanism is not relevant to our results. Moreover,
earlier it have been shown that the hysteresis of magnetoresistance is
independent on the in-plain component of magnetic field.Budantsev2009 It
allowed us to conclude that the hysteretic effect is not related to the
electron spin.
The measurements of non-equilibrium magnetization of a 2DEG also show stepwise
relaxation.Smith2011 It should be noted that transport measurements give a
valuable complementary information about the non-equilibrium state of a 2DEG.
Fact is that the magnetization is an integral characteristic determined by the
sum of all the magnetic moments, which are caused by all magnetization
currents in a sample, while the magnetoresistance is sensitive to a local
electrochemical potential of edge states. Particularly in this paper we
succeed to observe the correlation between the relaxation of magnetoresistance
of two remote constrictions and have established that NECs are localized near
the edge.
Figure 3: The dependence of the hysteresis loop area on the constriction width
at the filling factors $\nu$=1 and 2 at temperature 0.48 K.
To determine the width of spatial localization of NECs the dependence of
magnetoresistance hysteresis loop area on the lithographic width of the
constriction has been experimentally studied. For this purpose we used the
series of samples fabricated from the same heterostructure. The samples were
the Hall bars with constrictions of different widths. The constrictions of
widths from 0.8 $\mu$m to 1.3 $\mu$m have been fabricated by means of electron
beam lithography in the central part of Hall bars. All the constrictions have
low resistance in zero magnetic field (0.2$\div$0.8 k$\Omega$) and demonstrate
QHE in high magnetic fields. At the temperature of 0.48 K the hysteresis of
magnetoresistance have been observed at the filling factors $\nu$=1, 2 and 4
for all the samples. Fig. 3 shows the hysteresis loop area as a function of
the constriction width at the filling factors $\nu$=1 and 2. Experimental
points are well fitted by linear functions. One can see that the hysteresis
loop area decreases with the constriction width and vanishes at certain
critical width $W_{0}$. Moreover these critical widths coincide for both
filling factors and are about $W_{0}\approx$1.35$\mu$m. Assuming that the
critical lithographic width comprises of two counter-propagating NECs of the
width $W_{\mathrm{NEC}}$ and two depletion regions of the width
$W_{\mathrm{depl}}\approx$0.2$\mu$m:
$W_{0}=2W_{\mathrm{NEC}}+2W_{\mathrm{depl}}$, we obtain
$W_{\mathrm{NEC}}\approx$0.5$\mu$m is comparable with the depletion width. The
edge character of NECs has been also reported in the Ref. [Klaffs2004, ]
devoted to the study of spatial distribution of the non-equilibrium
electrostatic potential on a 2DEG surface, where it has been shown that the
most significant changes of potential take place close to the edge.
The observed hysteresis of magnetoresistance exhibits the following
properties: (i) it has two phases of relaxation, fast and slow; (ii) the slow
relaxation is occurred by jumps resembling the Barkhausen jumps in
ferromagnetics; (iii) it demonstrates inverted anti-coercive behavior; (iv) it
shows multistability and has an anhysteretic curve; (v) it can be observed in
sufficiently narrow channel thus pointing out its edge character.
## IV Discussion
The existing physical models are able to interpret only separate experiments.
Unfortunately there is no any model which would not contradict all the set of
experimental data on the non-equilibrium state in a 2DEG in the QHE regime.
Analyzing the experimental data obtained in our and other studiesBudantsev2007
; Klaffs2004 ; Matthews2004 ; Usher2009 ; Ruhe2009 we suggest such a
qualitative picture. We propose the picture of spatial distribution of NECs in
a 2DEG and discuss its relation to the experimentally observed non-equilibrium
magnetization of the 2DEG. For simplicity, we consider a 2DEG at the filling
factor near $\nu=1$, when electrons form the incompressible liquid on the
first Landau level and inter Landau level scattering is assumed to be
suppressed.
### IV.1 Magnetoresistance of constriction
The obtained experimental data on the magnetoresistance can be interpreted in
the frame of the following physical picture.
Down-sweep of the magnetic field gives rise to the vortex electric field that
induces outflow of the electrons from the bulk to the edge.Dolgopolov1992 ;
Laughlin1981 In this case the area of the incompressible liquid increases and
the edge channels are shifted closer to the lithographic borders of the
sample. It means that the counter-propagating edge currents in the
constriction become more distant from each other (Fig. 4(a-c)) than in the
case of equilibrium state and the backscattering is suppressed. In the type II
samples the magnetoresistance is indeed lower than its equilibrium value on
the anhysteretic curve. In the type I samples at filling factor $\nu=1$ the
decrease of magnetic field leads to complete suppression of the backscattering
in the constriction, manifested in the vanishing of the longitudinal
magnetoresistance and the establishment of the QHE regime.
Up-sweep of the magnetic field leads to vortex electric field, which in turn
induces outflow of the electrons from the edge to the bulk. In this case the
area of the incompressible liquid decreases and the edge channels are shifted
away from the lithographic border of the sample. The counter-propagating edge
currents in the constriction come closer together (Fig. 4(d-f)) than in the
case of equilibrium state and the backscattering increases. This results in
magnetoresistance rise as we have observed in both types of the samples. The
magnetoresistance of the type II samples is indeed higher than its equilibrium
value on the anhysteretic curve. In the type I samples at the filling factor
$\nu=1$ the increase of magnetic field switches the constriction from the QHE
regime to the resistive state. In other words, the magnetic field sweep leads
to magnetic field induced QHE breakdown.
Figure 4: Non-equilibrium electron redistribution in the macroscopic 2DEG (a),
(d) and in the constriction (b), (e) and top view of the 2DEG with narrowing
(c), (f) at up- and down-sweep of the magnetic field correspondingly at the
filling factor $\nu$=1. Bulk states in the constriction are occupied (c) or
not occupied (f) with incompressible liquid in different magnetic field sweep
direction. Arrows designate nonequilibrium currents.
### IV.2 Non-equilibrium magnetization
Let us consider how the proposed picture corresponds to the non-equilibrium
magnetization measurements. The spatial distribution of the NECs have been
discussed in Refs. [Usher2009, , Ruhe2009, , Matthews2004, ]. However, the
details of the mentioned models are not fully consistent with the experimental
data obtained up to date. In particular, the model proposed in Refs.
[Usher2009, , Matthews2004, ] assumes that the non-equilibrium charge is
redistributed over the whole plane of a 2DEG and consequently the electric
field is formed in the bulk. This contradicts experimental data obtained in
the present study and discussed in Ref. [Klaffs2004, ] pointing to the fact
that the non-equilibrium charge redistribution takes place only near the edge.
In Ref. [Ruhe2009, ] the spatial distribution of the magnetization currents at
the 2DEG edge is qualitatively described (Fig.8 in [Ruhe2009, ]), but the
provided picture corresponds rather to equilibrium state because the common
Fermi level for the bulk and the edge is introduced, up to which all the
electronic states are occupied. Besides in this case the observed
magnetization should not exceed the amplitude of the equilibrium dHvA
oscillations while the experimental non-equilibrium magnetization is about
20$\div$60 times larger.
Figure 5: The structure of the non-equilibrium edge currents at down- (a) and
up-sweep (b) of the magnetic field.
The decrease of magnetic field at the filling factor $\nu=1$ shifts the
electrons from the bulk to the edge within the first Landau level. Due to the
fact that the area occupied by the bulk states is much larger than that
occupied by the edge states one can neglect the changes of the electrostatic
potential in the bulk and consider that it remains flat. In this case the most
significant changes are occurred near the edge, where a dipolar strip
decreasing the electrostatic potential in the bulk and increasing it near the
edge is formed. The electrons on the first Landau level populate the edge
states up to the level of non-equilibrium electrochemical potential
$\mu_{down}$ exceeding the potential of the empty bulk states by the value
$\Delta\mu$ (Fig. 5(a)). In the crossed electric and magnetic fields the
electrons drift along the edge forming the magnetization current approximately
equal to $e^{2}/h\cdot\Delta\mu$. This current corresponds to the
experimentally observed non-equilibrium paramagnetic
magnetization.Matthews2004 ; Ruhe2009
The increase of magnetic field at the filling factor $\nu=1$ leads to more
complicated distribution of the magnetization currents. The outflow of the
electrons from the edge to the bulk leads to the formation of the dipolar
strip near the edge that rises the electrostatic potential in the bulk and
forms the potential well near the edge followed by the depletion region. It
means that the electrostatic potential is non-monotonic and has a minimum near
the edge, caused by interplay of two oppositely directed electric fields: the
field of overpopulated bulk and the depletion field (Fig. 5(b)). On either
side of the potential minimum the incompressible electron liquid forms
counter-propagating drifting currents — giant diamagnetic current flowing
closer to the bulk and low paramagnetic current flowing in close proximity to
the edge. The presence of the paramagnetic current is confirmed by the
experimental fact that the sign of the Hall resistance $R_{H}$ does not depend
on the magnetic field sweep direction in a given orientation of the magnetic
field. Hence the direction of the transport current, flowing along the edge of
a sample, coincides with the direction of the paramagnetic current as in case
of equilibrium state.
Edge states are populated up to the non-equilibrium electrochemical potential
level $\mu_{up}$ so that the empty edge states on the first Landau level are
lower than populated bulk states by the value of $\Delta\mu$. The total
magnetization current is dominated by the giant diamagnetic current
$e^{2}/h\cdot\Delta\mu$. To explain the experimentally observed giant non-
equilibrium magnetization of a 2DEG Usher2009 it is necessary that
$\Delta\mu\geq(20\div 60)\hbar\omega_{c}$, where $\omega_{c}$ is a cyclotron
frequency.
Here we do not discuss the finite magnetic field range of the hysteresis loop
as it requires detailed consideration of disorder in nonlinear screening
conditions.Chklovskii1992
### IV.3 Anti-coercive behavior
There are two possible reasons for anti-coercive behavior: (i) an abrupt
movement of the boundary of the incompressible liquid at small change of the
magnetic field, (ii) topological transitions in quantum Hall liquid. Discuss
them in more detail.
(i) at the filling factor $\nu=1$ the magnetic field change $\Delta B$ causes
the shift of the boundary of the incompressible liquid relative to the
lithographic edge. Besides, the larger the sample the smaller the change of
the magnetic field $\Delta B$ is required to shift this boundary by an amount
of the order of the depletion width. For example for disk-shaped 2DEG of
radius $R$:
$\Delta B\approx\frac{W_{depl}}{2R}\cdot B.$ (1)
(ii) As seen from Figs. 4(c) and 4(f) the shift of the boundary of the
incompressible liquid leads to topological transitions in quantum Hall liquid.
Fig. 4(c) corresponds to the total transmission of the edge channels through
the constriction, while Fig. 4(f) corresponds to the intense backscattering.
Thus, a small change of the magnetic field leads to a drastic change of the
magnetoresistance.
Topological transition has been observed in samples of type I: the vanishing
of the longitudinal magnetoresistance has been observed indicating total
suppression of backscattering in the constriction at down-sweep of the
magnetic field, while up-sweep of the magnetic field causes sharp jump of the
magnetoresistance from zero to 10 k$\Omega$. Abrupt movement of the boundary
of the incompressible electron liquid at small change of the magnetic field is
inevitably realized in the both types samples. This is manifested itself in
the sharp changes of the magnetoresistance in minor loops measurements.
Figure 6: Hysteresis of 2DEG magnetization (large loop) in QHE regime against
the de Haas - van Alphen oscillations (thick line) is inverted (anti-coercive
behavior) that is the response advances the magnetic field sweep. Small loop
demonstrates the retarding (coercive) behavior.
The discovered unusual anti-coercive behavior is in accordance with the law of
energy conservation. Consider the behavior of the non-equilibrium
magnetization of a 2DEG in the QHE regime. In Fig. 6 a typical hysteresis of
the magnetization on the background of dHvA oscillations at the filling factor
$\nu=1$ is shown. Matthews2004 The system response (magnetization) also
advances the external parameter (magnetic field) change as in the case
discussed. However the total work performed by the system $-\oint{Md\,B}$
(shadow area in Fig. 6) is negative, that is in accordance with energy
conservation law. Retarding response behavior would correspond to the positive
total system work (dashed area in Fig. 6) per cycle contrary to the law of
conservation of energy.
### IV.4 Comparison with ferromagnetics
Let us compare the observed hysteresis of magnetoresistance with the
hysteresis of magnetization of ferromagnetics.
It is well known that ferromagnetism is of spin origin. The exchange
interaction in ferromagnetics causes the spins align. At the same time, the
requirement of a minimum of the total energy of the system leads to its
splitting into domains with opposite spin orientations. A change of the
magnetic field results in spin flip process near the domain walls causing
their “motion”. However, due to a disorder the transformations of domain
structure is retarded with respect to the external magnetic field change. This
is the reason of hysteresis with positive coercivity in ferromagnetics. While
in a 2DEG a small changes of the external magnetic field lead to sharp motion
of the boundary of the incompressible liquid and causes the hysteresis with
negative coercivity.
Earlier it has been found Budantsev2007 that the hysteresis of
magnetoresistance is suppressed by temperature. It has been shown that the
amplitude of hysteresis linearly increases as the temperature decreases and
saturates at certain low temperature ($\sim$400 mK). The non-equilibrium
magnetization of a 2DEG has similar temperature dependence.Matthews2004 The
hysteresis of magnetization of ferromagnetic materials is also suppressed by
temperature but it has different functional dependence described by Weiss
theory of ferromagnetism. Bertotti1998
It should be noted that NECs in the QHE regime are induced by sweep of the
magnetic field and can be considered as eddy currents as they are often
called. However, in ferromagnetics it is enough to reduce the sweep rate to
eliminate the influence of eddy currents. In this sense eddy currents in
ferromagnetics are time-dependent because of finite longitudinal conductivity
$\sigma_{xx}$. While in a 2DEG in the QHE regime when $\sigma_{xx}\approx 0$
the generation of NECs is inescapable at any reasonable sweep rate. According
to our results NECs relaxation lasts for hours after stopping the magnetic
field sweep. In Ref. [Kershaw2007, ] NECs have been shown to circulate for
many hours. In terms of Ref. [Bertotti1998, ] under conventional experimental
conditions NECs can be considered as time-independent.
## V Conclusions
We have shown that sweep of the magnetic field in the QHE regime results in
drastic changes of electrostatic potential in the bulk with respect to the
edge if the inter Landau levels scattering is suppressed. This leads to giant
magnetization currents flowing along the edge and topological transitions in a
2DEG with a constriction. The latter results in changes of magnetoresistance.
The hysteresis of magnetoresistance of a 2DEG in the QHE regime has
significant phenomenological similarities with the hysteresis of magnetization
of ferromagnetic materials, showing multistability, jumps of relaxation,
temperature suppression and having the anhysteretic curve. At the same time a
fundamental difference expressed in anti-coercive hysteresis behavior of the
magnetoresistance has been found.
The possible causes of the limitations of the inter Landau level scattering as
well as the possible limiting mechanisms of NECs are beyond the scope of this
article and require further investigation. However we have experimentally
shown that NECs are induced in a narrow ($\sim$0.5$\mu$m) area along the edge.
Despite the large number of studies devoted to the non-equilibrium state of a
2DEG induced by the sweep of the magnetic field in the QHE regime there is no
clear understanding of the phenomena at a microscopic level as well as a
systematic phenomenological study of the dependence of the phenomena on
parameters of heterostructures and conditions of epitaxial growth up to date.
Unfortunately the NECs are insensible in a conventional transport measurement
in a linear response regime due to zero longitudinal resistance. In such
measurements the condition of the 2DEG is not clear in advance. Such
uncertainty can bring unexpected difficulties in studies of a 2DEG in the QHE
regime, such as the study of QHE breakdown by the electric current, that has
no full and consistent model up to date (see for a review Ref. [Nachtwei1999,
]). The present study raises the issue of the need to search critical
parameters that define the phenomena and explain the phenomena at a
microscopic level. Therefore, some existing concepts on the phenomena in 2DEG
in the QHE regime can be extended or revised taking into account the found
non-equilibrium.
## Acknowledgments
The reported study was supported by RFBR (research project
No.14-02-31740-mol-a and No.12-02-00532-a), Program of fundamental scientific
research DNIT RAS (project 3.2).
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|
arxiv-papers
| 2013-09-24T09:09:46 |
2024-09-04T02:49:51.423544
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "M. V. Budantsev, D. A. Pokhabov, A. G. Pogosov, E. Yu. Zhdanov, A. K.\n Bakarov, and A. I. Toropov",
"submitter": "Dmitry Pokhabov A",
"url": "https://arxiv.org/abs/1309.6089"
}
|
1309.6117
|
# Efficient Gate-tunable light-emitting device made of defective boron nitride
nanotubes: from ultraviolet to the visible
Claudio Attaccalite1, Ludger Wirtz2,3, Andrea Marini4, Angel Rubio5,6111To
whom correspondence should be addressed; email: [email protected] 1 Institut
Néel, CNRS, 25 rue des Martyrs BP 166, 38042 Grenoble cedex 9 France
2 Institute for Electronics, Microelectronics, and Nanotechnology (IEMN), CNRS
UMR 8520, Dept. ISEN, 59652 Villeneuve d’Ascq Cedex, France
3 Physics and Material Sciences Research Unit, University of Luxembourg, 162a
ave. de la Faïencerie, L-1511 Luxembourg
4 Istituto di Struttura della Materia (ISM), Consiglio Nazionale delle
Ricerche, Via Salaria Km 29.5, CP 10, 00016 Monterotondo Stazione, Italy
5 Nano-Bio Spectroscopy Group and ETSF Scientific Development Centre,
Departamento de Física de Materiales, Centro de Física de Materiales CSIC-
UPV/EHU-MPC and DIPC, Universidad del País Vasco UPV/EHU, Av. Tolosa 72,
E-20018 San Sebastián, Spain
6 Fritz Haber Institut der Max Planck Gesellschaft, Faradayweg 4-6, 14195,
Berlin,Germany
Boron nitride is a promising material for nanotechnology applications due to
its two-dimensional graphene-like insulating and highly-resistant
structure.Pakdel2012256 ; wirtz_review ; RevModPhys.82.1843 Recently it has
received a lot of attention as a substrate to grow and isolate
graphenebngraphene as well as for its intrinsic UV lasing response.watanabe1
; watanabe2 Similar to carbon, one-dimensional boron nitride nanotubes
(BNNTs) have been theoretically predictedrubiobn and later
synthesised.tubesynthesis Here we use first principles simulations to
unambiguously demonstrate that i) BN nanotubes inherit the highly efficient UV
luminescence of hexagonal BN; ii) the application of an external perpendicular
field closes the electronic gap keeping the UV lasing with lower yield; iii)
defects in BNNTS are responsible for tunable light emission from the UV to the
visible controlled by an transverse electric field (TEF). Our present findings
pave the road towards optoelectronic applications of BN-nanotube-based devices
that are simple to implement because they do not require any special doping or
complex growth.
Scientists have worked hard in the last decades to grow defect free nano-
structures. The near-perfect atomic arrangement at the nano-scale has been
employed to create new efficient devices as light-emitters, transistors and
sensors. While many applications benefit from using defect free materials, the
presence of particular impurities can generate new and fascinating properties.
For example F-centre in ionic crystals have been widely used in luminescence
applications.RevModPhys.18.384 More recently nitrogen vacancies in diamond
have been proposed for quantum computation.PhysRevA.71.060310 Defects play an
important role also for the luminescence properties of hexagonal boron nitride
and related nanostructures.PhysRevLett.100.189701 ; attaccalite2011 ;
PhysRevB.75.085205 ; PhysRevB.77.235422 ; PhysRevB.78.155204
Figure 1: Panel $a)$ left(right): lowest conduction band (highest valence
band) in presence of an external field E=0.206 V/Å. In panel $b)$ we report
the band gap of BN $8\times 8$ nanotube versus the TEF in LDA and $G_{0}W_{0}$
approximations; in panel $c)$ the corresponding optical absorption computed at
the many-body Bethe-Salpeter level.
Similarly to graphene, a single BN layer can be rolled up to form new
structures ranging from single and multi-wall nanotubes to BN-
fullerenswirtz_review . In contrast to graphite, the ionic character of the BN
bond results in a wide band-gap of about 6 eV for bulk hexagonal-
BNPhysRevLett.96.126104 ; PhysRevLett.100.189701 ; arnoud2006 . The
combination of a such large gap with a strong electron-hole attraction makes
the optical properties of hexagonal-BN based nanostructures largely
independent of the layer arrangement and dimensionality.PhysRevLett.96.126104
; wirtz_review Although bulk h-BN has been shown to exhibit a strong
luminescence, it cannot be used for optical applications in the visible range
because the emission frequency is fixed to about 5.75 eV in the UV. However
the presence of impurities can drastically modify this scenario, as it has
been shown, theoreticallyattaccalite2011 and
experimentally.PhysRevB.75.085205 ; PhysRevB.78.155204 ; PhysRevB.77.235422 .
The defects electronic structure and formation energies have been widely
studied in bulk h-BN and BN nanotubes.PhysRevB.67.113407 ; p911623m ;
jp056941l ; berzina ; jp800096s ; zobelli Luminescence in the visible was
attributed to the presence of deep levels in the sample,ADMA200700493 ;
nl0726151 ; Chen2010S353 whereas the UV emission is an intrinsic response
(Frenkel-exciton) of the structurally perfect hexagonal
BN.pierret2013excitonic Here we propose to engineer BNNTs with particular
defects in order to generate light-emission in a wide range of frequencies,
that can be tuned by means of an external electric field. The range of
tunability of the proposed nanotube-based light emitting device depends from
defect location and type.
For pure BNNTs, it has been shown that the application of a transverse
electric field generates a Stark effect leading to a strong reduction of the
band gapPhysRevB.69.201401 ; PhysRevLett.94.056804 . The external field leads
to a localisation of the conduction-band minimum/valence-band maximum on
opposite sides of the tube(see Fig. 1). The corresponding energy shift of the
band edges is thus proportional to the nanotube diameter. Experimentally as-
grown nanotubes contain defects that lead to both deep or shallow levels in
the gap. The wave-functions of these levels are, a priori, only slightly
affected by the presence of an external electric fieldattaccalite2007 ;
jp800096s because they are associated with localised orbitals centred on the
impurity. However, their level position with respect to the bands edges
changesjp800096s ; attaccalite2007 because valence(conduction) bands are
modified by the external field. Here we will show that this property can be
employed to produce tunable and highly efficient bright light-emission devices
based on defective BN nanotubes.
## Results
We model the electronic and structural properties of the pure and defected BN
nanotubes under a TEF using state-of-the-art first-principles methods based on
density functional theory(DFT) combined with many-body perturbation-
theory(MBPT) approaches. These methods allow calculation of quasiparticle-band
structure and optical properties with a high degree of accuracy (see Methods
section for details). In the past this theoretical framework has been shown to
be very efficient in predicting the electronic properties on BN nanotubes that
were later on confirmed in the experiments.wirtz_review ; rubiobn ;
PhysRevLett.96.126104 ; RevModPhys.82.1843
We start our study by analysing the case of pure isolated BN nanotubes
immersed in a static transverse electric field. A transverse electric field
reduces the band gap, as shown in panel $b)$ of Fig. 1. The gap reduction,
induced by the TEF, is directly proportional to the electric field strength
and to the tube diameter.PhysRevB.78.085423 Surprisingly the shrinking of the
band-gap slightly modifies the optical response of the tube.attaccalite2007
The main exciton remains in the same position, while a small fraction of its
spectral weight is redistributed to higher excitons (see panel $c)$ in Fig.
1). In fact the conduction and valence orbitals contributing to the gap
reduction are localised on the opposite side of the tube, they have a very
little overlap (see panel $a)$ of Fig. 1) and therefore their contribution to
the optical response is negligible. Therefore we conclude that while the giant
Stark effect present in BNNTs can modify their transport propertiestransport ,
it leaves the optical response mainly unchanged. The light emission spectra of
pure BNNTs is thus not tunable by an external electric field. The presence of
defects drastically modifies this picture. Different experiments have shown
that impurities induce light emission below 5 eV in BN
nanostructures,PhysRevB.75.085205 ; PhysRevB.78.155204 ; PhysRevB.77.235422
and modify the luminescence arising from the main bulk exciton.attaccalite2011
These effects can be explained by the presence of deep levels in the BN band
gap.attaccalite2011 The low frequency emissions are due to transitions from
and to these levels. Moreover when the impurity levels are close to the top
valence band or bottom conduction band they mix with the bulk excitations
giving rise to a splitting of the main excitonic peak.attaccalite2011 Among
the different impurities responsible for light emission we can distinguish two
families: defect complexes and single defect centres. The first family is
formed by multiple defects as for instance di-vacancies, defect lines and so
on. The second family consists in a single defect centre as for instance boron
(nitrogen) vacancies or a substitution of a boron(nitrogen) atom with a carbon
one. The main difference between these two families lies in the different kind
of transitions involved in light absorption/emission processes. In the case of
defect complexes both donor and acceptor states are present in the band gap
while in the other case there is only a single donor or acceptor state.
Therefore in the case of defect complexes the optical response is dictated by
the quasi-donor-acceptor transitionsMuseur2008 , while in the simple defect
centres light absorption/emission is due to transition between bulk states and
deep defect levels (see Fig. 2).
Figure 2: In the top panel we report the oscillator strength versus the angle
between the defect and the bottom conduction band for $V_{B}$. The optical
matrix elements are averaged on the first conduction bands within an energy
range of $0.15$ eV. The straight lines are the dipole matrix elements at zero
TEF. In the same figure it is present also a schematic representation of the a
BN tube with a defect in presence of a TEF. In the bottom panels we show the
band structure of the same tube versus the angle between the defect and the
electric field. The red arrow represents the transition responsible for the
luminescence in presence of $V_{B}$.attaccalite2011 The intensity of the TEF
is 0.206 V/Å.
The electronic structure of defects in BN nanotubes is similar to the one of
defects in a single BN-sheet.PhysRevB.67.113407 In fact due to the large band
gap curvature effects play a minor role on the optical properties of pure
BNNTs, where the strong localisation of excitons renders the optical spectra
almost independent from the tube diameter and chirality.louietube ;
PhysRevLett.96.126104 In order to simulate a tube with defects we use the
same methodology employed for the pure tube but with larger supercells in such
a way to reduce the defect-defect interaction. Although large part of the
tubes produced in the experiments are multi-wall and possess a zig-zag
chirality, we chose a $12\times 12$ armchair one as prototype for our study.
This choice is motivated by two reasons: first the primitive cell of an
armchair tube, radius being equal, contains less atoms than a chiral one,
second we expect only small differences with respect to the optical response
of multi-wall or chiral nanotubes for the reasons discussed above. When we
turn on a TEF, the band gap of the tube shrinks and consequently, the defect
levels change position with respect to the band edges.jp800096s ;
attaccalite2007 The orbitals associated with defects levels are strongly
localised on the impurities (see right panel of Fig. 3) and therefore they are
slightly deformed by the presence of the external field. To first order, the
shift of the defects levels is thus given by the potential generated by the
TEF and depends therefore on the position of the defect respect to the
direction of the electric field (see inset in Fig. 2). This is visualised in
the bottom panel of Fig. 2 for three different defect positions. We will show
in the following, how this property gives rise to a tunable and efficient
light emission.
In order to predict the emission frequencies of BN nanotubes in presence of
defects, we used a simplified approach. The first necessary ingredient to get
light emission is non-vanishing optical matrix elements between the discrete
donor(acceptor) state and the continuum states of the bottom conduction(top
valence) bands. In the upper panel of Fig. 2 we show the strength of the
optical matrix elements between the defect level and the bottom conduction
bands for the case of a Boron vacancy, $V_{B}$. The optical matrix element
displays a strong dependence on the angle between the defect position and the
electric field (see inset in the top panel of Fig. 2). The same phenomena
occurs for other simple acceptor or donor defects like substitution of a
Nitrogen or a Boron atom by a carbon one, $C_{B}$ and $C_{N}$ respectively.
Furthermore, we note that the optical matrix elements for polarisation along
the tube axis (z-axis) dominate, which also holds for the optical response of
pure nanotubeswirtz_review .
The presence of the external electric field localises valence and conduction
bands on opposite sides of the tube (see Fig. 1), therefore transitions are
maximal only when the defect is aligned with the bottom(top) of the
conduction(valence) band. The dipole element decreases to almost zero as the
defect is turned to the opposite side of the tube. Consequently one can expect
that luminescence, which is generated by transitions from and to the defects
levels, will be efficient only when the defects are positioned on the side of
”localised” conduction(valence) band edge. This focusing effect increases with
the tube size and field intensity.
Now that we are sure that transitions from and to simple defect centres are
not zero in presence of a transverse electric field, we investigate how the
field modifies the emission frequencies. In order to predict light emission we
start from the quasi-particle(QP) band structure in presence of defects. We
consider the energy differences for transitions between defect states and the
top valence(bottom conduction) states. This allows us to investigate light-
emission versus transverse electric field, without including electron-hole
interaction or lattice relaxation (see Methods section). We found that also in
presence of defects the GW renormalisation for conduction(valence) bands and
defects levels is almost a constant respect to the external electric field
(see also Fig. 1(b)).
---
Figure 3: Schematic representation of light emission process for an acceptor
impurity in a BN nanotube. On the left a simplified band structure picture in
presence of a TEF. On the right conduction and defect orbitals responsible for
the emission process.
Luminescence can be estimated from the QP band structure in presence of TEF as
the sum of independent transitions between conduction bands and defect states.
However light emission originates from electron-hole recombination, a two-
particle process that cannot be described by means of the quasi-particle band
structure only. In fact electrons and holes attract each other and this
attraction modifies the transition energies. These processes can be naturally
treated within a two-particle Green’s functions formalismstrinati and it has
been shown that transition energies from and to defect states are strongly
renormalisedattaccalite2011 by the electron-electron correlation. In order to
model this correction, we calculate the exchange and electron-hole attraction
between the defect level and the bottom conduction(top valence) bands only. In
the past this approximation has been successfully employed to predict
excitation energies of F-colour centres.louie
In addition to the corrections originating from the electronic correlation, we
have to consider the contributions due to the lattice relaxation induced by
the excited carriers. These are the so called Stokes and anti-Stokes shifts.
The Stokes shifts can be estimated by means of a constrained DFT calculation
with different defect occupation. We investigated three different defects, an
acceptor the Boron vacancy $V_{B}$, a donor the Carbon substitution of a Boron
atom $C_{B}$, and the Boron-Nitrogen di-vacancy $V_{BN}$. In order to estimate
the Stokes shift we considered the case of a completely empty acceptor
state(or a completely filled donor state). In principle one should consider
also the correction coming from the partial filling (emptying) of the
conduction(valence) bands, but this is supposed to be a minor effect because
these bands are delocalised along the $\vec{z}$ direction. In this way we
obtain a rough estimation of the Stokes shift of $\Delta E_{s}\simeq
0.19~{}eV$ for $C_{B}$ and $\Delta E_{s}\simeq 0.03~{}eV$ for $V_{B}$. We did
not calculated any Stokes shift for $V_{BN}$ because in this case it is
irrelevant as it will be clear in the following. Vested with this theoretical
approach we proceed in the study of light emission versus the external
electric field.
Figure 4: Predicted light emission for different defects as a function of the
transverse electric field. All the defects are taken in the position of
maximum emission, according to their optical matrix elements, see also Fig. 2.
In the $V_{BN}$ we did not include any Stokes shift.
We report our predicted light emission for $BN(12,12)$ tube versus the
transverse electric field in Fig. 4. As one can see from the figure an
external electric field allows to vary the emission frequency in a large
spectral range for the CB and VB cases. Notice that in presence of defect
complexes, as for instance the BN di-vacancy $V_{BN}$, the emission frequency
does not change with the external field. In fact in this case the emission is
dominated by transitions between donor and acceptor states in the band-
gap.Museur2008 ; attaccalite2011 Since the wave-functions associated to these
states are localised on the impurity, the effect of external electric field is
irrelevant. In the left panel of Fig. 3 we present a schematic representation
of the light emission process from BN nanotubes in presence of defects. We
want to underline that this process happens only when the defect is aligned
with the conduction(valence) maximum, otherwise the emission will be
inefficient due to the small dipole matrix elements.
Although the results of Fig. 4 can be theoretically extended to larger tubes,
calculations become soon prohibitive due to the large number of atoms, the
vacuum in the super-cell and the number of conduction bands that enter in the
many-body operators. Therefore in order to predict light emission in larger
(more realistic) tubes we assume many-body corrections to be a constant with
respect to the tube size and we fit the emission energy with a simple linear
curve
$E_{emission}=E_{0}+\alpha\xi$ (1)
where $\xi$ is the external electric field. This relation was already employed
to describe the band gap closing of h-BN nanotubes under the effect of a TEF
in simple tight-binding models and _ab-initio_ calculations.PhysRevB.78.085423
; Chegel2012154 In principle the linear coefficient $\alpha$ depends on the
tube size. In order to estimate this dependence we performed different
calculations at the DFT level, varying the tube size. We found that $\alpha$
changes linearly with the tube radius $R$, $\alpha(R)=\alpha_{0}+R\beta$. A
similar behaviour has been found for the pure BNNTs gap versus the electric-
field and tube radius.PhysRevB.78.085423 Combining the previous two equations
we can predict the electric field $\xi$ necessary to produce light emission at
a given frequency $E_{1}$:
$\xi=\frac{E_{1}-E_{0}}{\alpha_{0}+R\beta}.$ (2)
Now we use Eq. 2 to estimate the intensity of the TEF that will induce
emission in the visible range ($1.65-3.1$ eV). In Fig. 5 we report the visible
emission range for the $V_{B}$ and $C_{B}$ cases versus the TEF intensity and
tube radius.
Figure 5: Visible light emission range as function of the tube radius and the
external electric field for $V_{B}$ and $C_{B}$ cases (rainbow colours are
just a guide to the eyes).
In general an increase of the tube size reduces the strength of the transverse
electric field necessary to obtain emission in the visible range, For
sufficient large tubes the TEF intensity is of the same order of the one
available in small devices. Notice that an TEF produces an electrostatic
potential inside the tube that is proportional to the TEF intensity and the
tube radius.jp800096s Therefore Eq. 2 breaks down for too large nanotubes or
too strong fields. However a giant Stark effects has been experimentally
measured in pure BN nanotubes with a radius of about $23~{}\AA$ and a TEF of
$0.08~{}V/\AA$. In the same experiment a gap reduction of more than $1~{}eV$
has been obtained.PhysRevLett.94.056804 Comparing these values with our
extrapolation in Fig. 5 it is clear that there is a large margin to produce
visible light with experimentally accessible nanotubes and electric fields.
Finally we consider defects formation and their charge state. Recent
experimentsjin2009fabrication have shown that it is possible to introduce
defects in h-BN structures by means of electron irradiation. This process is
mainly dominated by boron mono-vacancies even if other larger vacancies are
present. These vacancies can also be transformed in substitutional defects by
introducing C atoms in the experiment,krivanek2010atom ; risto and the final
process can be controlled by charging the system during the irradiation.risto
These advances make possible the realisation of the device that we are going
to discuss in the following. Regarding the charge state of the defects, in the
present paper we investigated only neutral ones. Charged defects posses
different relaxation energies and electronic structure. This fact influences
also their optical properties, as it has been recently shown in the case of
vacancies in SiC.bockstedte2010many The present results can easily be
extended to charged defects and we expect that the main findings will remain
valid. In fact the tunability of the light emission it is related to the
localisation of defect states versus the delocalised bulk ones. Therefore a
different charge state will modify the emission at zero field but not its
behaviour in presence of a TEF.
Now that we have shown how to produce tunable light emission with defective BN
nanotubes and discussed the feasibility of our idea, we briefly present the
possible configurations of a device based on BN nanotubes. The generic
configuration of the device (see Fig. 6) comprises depositing as-grown BN
nanotubes on an insulating surface (for example silicon oxide) acting as a
dielectric to enable the application of the gated electric field that controls
the light emission. The configurations is very much similar to the one of a
field effect transistor (FET). The activation of the BN-defected
optoelectronic device could be done by one of the following three processes:
i) using UV light, ii) introducing an ambipolar current that recombines in the
defect and emits light dictated by the applied gate voltageChen18112005 iii)
using tunnelling current through an STM tip close to the nanotube. The excited
electrons would inelastically decay very fast into the lowest energy state
(the defect-liked Frenkel exciton) that would further decay by emitting light,
again with a frequency dictated by the applied voltage, a process similar to
the one leading to light emission in electronically excited semiconductors and
fluorescent materials. A schematic set-up of those devices is illustrated in
Fig. 6.
## Discussion
In conclusion, we have shown that light emission from BNNTs with simple defect
centres can be tuned by the presence of TEF. This opens the possibility to use
these systems as light emitting devices. The use of (non-tunable) UV-light
emitting devices based on crystalline hexagonal BN has already been suggested
beforepatent . Here, we move one important step further by showing how
nanotube based devices could operate in the UV and visible range by varying
the external field. The external electric field, necessary to tune the
emission, can be applied using a field-effect transistor
configurationtransport . The intensity necessary to produce visible light
decreases with the tube size. The present results have been patented by some
of the authors.tubepatent
Figure 6: Schematic set-up for the suggested three possibilities to activate
the optoelectronic device based on defective BN nanotubes (i) light-induced
luminescence (ii) ambipolar transistor configuration (iii) electron induced
electron-hole pairs by means of STM tip.
Finally we envision that the present findings can be applied to other two-
dimensional semiconducting or insulating layered materials that form tubes, as
it is the case for transition metal dichalcogenides.cogenides
## Methods
BNNTs are simulated by using a supercell approach,rubiobn where the tube is
oriented along the $z$ direction, and a large empty space is left in the other
two directions between tube replica in order to reduce the tube-tube
interactions. Subsequently a sawtooth electric field (see inset in Fig. 1(a))
with the cell periodicity is added along the $x$ direction. In order to
simulate light emission in BNNTs we employed a combination of Density-
Functional Theory (DFT) plus Many Body Perturbation Theory (MBPT). DFT is an
exact theory for ground state properties and it is known to describe very well
the structural properties of boron-nitride nanostructures within Local Density
Approximation (LDA). All DFT calculations have been performed using a $1\times
1\times 5$ supercell containing 240 atoms. The distance between the tube
replica was 29 a.u. , we used a $1\times 1\times 2$ k-point sampling, LDA for
the exchange correlation functionalceperley , a plane waves cutoff of 45 Ry
for the wave-function and norm-conserving pseudo-potentials.troullier All DFT
calculations have been performed with the PWSCF codepwscf and the atomic
structures have been relaxed using a BFGS quasi-Newton algorithm. Excited
state and optical properties have been studied by means of MBPT. We calculated
quasi-particle properties solving a Dyson equation within the so-called
$G_{0}W_{0}$ approximationaryasetiawan1998gw ; gw-strinati , where all the
Green’s functions and the self-energy operator are constructed with
eigenvalues and eigenvectors of the Kohn-Sham(KS) Hamiltonian. Non-self
consistent GW calculations have been performed with the code YAMBOyambo using
a plasmon pole approximation for the dielectric constant. We used 30.000
G-vectors for the wave-function, 2 Ha for the response block size and 3000
bands for the screening. A cylindrical cutoff has been applied to the Coulomb
potential in order to reduce the tube-tube interaction. Neutral excitations,
responsible for the absorption spectra were obtained from a two-particle
Green’s function equation, the Bethe-Salpeter equation, that is solved in the
static ladder approximationstrinati , including excitonic effects. We excluded
quasi-free electron statesPhysRevB.69.201401 in the Bethe-Salpeter equation,
because they are not supposed to be responsible for luminescence. We performed
all calculations without including spin-polarisation effects. Even if we know
that exchange-splitting slightly modifies the defects levels
positionsattaccalite2011 ; PhysRevB.76.014405 , this effect does not modify
the main results of the paper. For the large tubes employed to get the results
in Fig. 4 we estimated the GW and electron-hole interaction from the one of a
BN-sheet with the same defects and a distance between the periodic replica
equal to the inter-tube distance.attaccalite2011
## Acknowledgements
We acknowledge financial support from the European Research Council Advanced
Grant DYNamo (ERC-2010-AdG - 267374) Spanish Grant FIS2011-65702-C02-01, Grupo
Consolidado UPV/EHU del Gobierno Vasco (IT-319-07) and European Commission
project CRONOS (280879-2). Computational time was granted by i2basque and BSC
”Red Espanola de Supercomputacion” and GENCI-IDRIS (Nos. 100063 and No.
091827). A. M. acknowledges funding by MIUR FIRB Grant No. RBFR12SW0J.
## Author Contributions
CA, LW, AM and AR contributed to the discussions, theoretical analysis and
writing of the manuscript. CA performed the calculations and AR designed the
research.
## Author Information
The authors declare no competing financial interest. Correspondence and
request or materials should be addressed to AR ([email protected]).
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|
arxiv-papers
| 2013-09-24T11:38:06 |
2024-09-04T02:49:51.433552
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Claudio Attaccalite, Ludger Wirtz, Andrea Marini and Angel Rubio",
"submitter": "Claudio Attaccalite CA",
"url": "https://arxiv.org/abs/1309.6117"
}
|
1309.6176
|
# Feature Learning with Gaussian Restricted Boltzmann Machine for
Robust Speech Recognition
###### Abstract
In this paper, we first present a new variant of Gaussian restricted Boltzmann
machine (GRBM) called multivariate Gaussian restricted Boltzmann machine
(MGRBM), with its definition and learning algorithm. Then we propose using a
learned GRBM or MGRBM to extract better features for robust speech
recognition. Our experiments on Aurora2 show that both GRBM-extracted and
MGRBM-extracted feature performs much better than Mel-frequency cepstral
coefficient (MFCC) with either HMM-GMM or hybrid HMM-deep neural network (DNN)
acoustic model, and MGRBM-extracted feature is slightly better.
Index Terms: restricted Boltzmann machine, robust speech recognition, feature
learning
## 1 Introduction
Since hybrid hidden markov model (HMM)-deep neural network (DNN) was
introduced to large-vocabulary continuous speech recognition (LVCSR), the
accuracy of speech recognition system has made significant performance
improvements in idealized environments [1][2]. Such progression urges the
development of speech recognition systems that are robust to background noise
and channel distortion, as more and more speech applications are deploying on
mobile devices. State-of-the-art robust automatic speech recognition system
usually involves intensive specialized domain knowledge [3]. But we are more
interested in the transformation of feature.
Feature learning (representation learning) [4] is a developing field that
grows alongside with deep learning. The aim of feature learning is to learn a
certain kind of transformation through which we are able to extract
information that makes discrimination much easier for classifiers. Feature
learning needs as little feature engineering as possible, and transformed
feature is much closer to real underlying factors that generate the original
features that we observe.
In this paper, we made our first attempt to apply the idea of feature learning
to robust speech recognition. Although dozens of alternatives have been
proposed over the past few decades, MFCC is still the default choice of
feature for many speech applications. Instead of trying to propose another
alternative, we are more interested in learning a better representation of
MFCC feature with restricted Boltzmann machine (RBM) and its variants. We will
explain why RBM may be able to learn a more suitable representation of MFCC
for robust speech recognition. We will also be proposing a new variant of RBM
called multivariate Gaussian restricted Boltzmann machine (MGRBM) which is
specially designed for modeling the distribution of speech data. MGRBM is able
to capture the evolving characteristic of speech within a context of several
frames which is difficult to model with a Gaussian restricted Boltzmann
machine (GRBM). We perform our experiments on the Aurora2 corpus and our
results show that the learned features are better than the original feature
for robust speech recognition.
## 2 Model
### 2.1 Restricted Boltzmann Machine
The Boltzman machine is a special kind of Markov random field which models the
joint probability distribution of the visible variable and hidden variable.
Visible and hidden variable are both defined to be multidimensional Bernoulli
variables. The distribution can be written as:
$\displaystyle p(v,h)=\frac{1}{Z}e^{-E(v,h;\theta)}$ (1)
and
$\displaystyle
E(v,h;\theta)=-\frac{1}{2}v^{T}Uv-\frac{1}{2}h^{T}Vh-v^{T}Wh-a^{T}v-b^{T}h$
(2)
$E(v,h;\theta)$ in (1) is called energy function. $\theta=\\{U,V,W,a,b\\}$.
$U,V,W$ models the visible-visible, hidden-hidden, and visible-hidden
interaction respectively. $a$ and $b$ are bias vectors.
The Restricted Boltzmann Machine (RBM) [5] is perhaps the most widely-used
variant of Boltzmann machine. The energy function of RBM is the simplified
version of that in the Boltzmann machine by making $U=0$ and $V=0$. That is,
the energy function of an RBM is:
$\displaystyle E(v,h;\theta)=-a^{T}v-b^{T}h-v^{T}Wh$ (3)
An RBM is typically trained with maximum likelihood estimation. Taking the
derivative with respect to the logarithm of the product of all the probability
of training cases, we can derive the learning algorithm of RBM as follows:
$\displaystyle\Delta W$ $\displaystyle=$
$\displaystyle\epsilon(<vh^{T}>_{data}-<vh^{T}>_{model})$ (4)
$\displaystyle\Delta a$ $\displaystyle=$
$\displaystyle\epsilon(<v>_{data}-<v>_{model})$ (5) $\displaystyle\Delta b$
$\displaystyle=$ $\displaystyle\epsilon(<h>_{data}-<h>_{model})$ (6)
The symbol $<\cdot>_{data}$ in (4)(5)(6) represents an average with respect to
the conditional distribution $p(v|h)$ and $<\cdot>_{model}$ represents an
average with respect to the joint distribution $p(v,h)$. The $<\cdot>_{data}$
for RBM is generally easy to train because:
$\displaystyle p(h_{j}=1|v)=sigmoid(b_{j}+W_{j}v)$ (7) $\displaystyle
p(v_{i}=1|h)=sigmoid(a_{i}+W_{\cdot i}^{T}h)$ (8)
(7)(8) can be derived from the definition of RBM and
$sigmoid(x)=\frac{1}{1+e^{-x}}$. However, $<\cdot>_{model}$ is much harder to
obtain. To address this problem, $<\cdot>_{model}$ is usually approximated
with $<\cdot>_{recon}$ as the following:
$\displaystyle\Delta W$ $\displaystyle=$
$\displaystyle\epsilon(<vh^{T}>_{data}-<vh^{T}>_{recon})$ (9)
$\displaystyle\Delta a$ $\displaystyle=$
$\displaystyle\epsilon(<v>_{data}-<v>_{recon})$ (10) $\displaystyle\Delta b$
$\displaystyle=$ $\displaystyle\epsilon(<h>_{data}-<h>_{recon})$ (11)
The $<\cdot>_{recon}$ represents an average with respect to the reconstruction
of the visible data. The reconstruction of visible data is obtained by setting
each node in hidden layer value 1 with probability (7), followed by setting
each node in visible layer value 1 with probability (8). This is the
contrastive divergence algorithm (CD) [6] for training of RBM. CD has been
empirically showed to be adequate for many applications.
### 2.2 Gaussian Restricted Boltzmann Machine
To model real-valued data, the Gaussian restricted Boltzmann machine (GRBM)
has been proposed [7][8]. The energy function of GRBM is typically defined
with:
$\displaystyle
E(v,h;\theta)=\sum_{i}\frac{(v_{i}-a_{i})^{2}}{2\sigma_{i}^{2}}-\sum_{ij}W_{ij}\frac{v_{i}}{\sigma_{i}}h_{j}-\sum_{j}b_{j}h_{j}$
(12)
in which $\theta=\\{a,b,W,\sigma\\}$ and $\sigma_{i}$ models the standard
deviation of each visible units.
Conveniently, the learning algorithm of GRBM is the same with RBM (9)(10)(11).
The conditional probabilities necessary for CD of a GRBM are:
$\displaystyle
p(v_{i}|h)=\mathcal{N}(a_{i}+\sigma_{i}\sum_{j}W_{ij}h_{j},\sigma_{i}^{2})$
(13) $\displaystyle
p(h_{j}=1|v)=sigmoid(b_{j}+\sum_{i}W_{ij}\frac{v_{i}}{\sigma_{i}})$ (14)
Generally speaking, $\sigma_{i}$ can be learned from data, but it’s difficult
with CD. The training data to be modeled with a GRBM are always pre-processed
mean 0 and variance 1 and thus $\sigma_{i}$ can be fixed with 1 and not
trained. The reason why $\sigma_{i}$ is difficult to train with CD can be
explained as follows: When $\sigma_{i}$ is much smaller than 1, the visible-
hidden effects (13) tends to be large and hidden-visible effects (14) tends to
be small. The result of such effect is that hidden units always tend to be
firmly 1 or 0, and thus undermine the whole training process.
One disadvantage of GRBM is its conditional independence assumption. That is,
conditioned on the hidden layer, each visible unit is assumed to follow a
Gaussian distribution and independent with each other. However, for natural
data such as speech and image, they tend to have local similarity property.
Take speech data for example, the smoothness of speech always makes one frame
of acoustic feature similar to the frames next to it. Such local similarity
property is difficult to capture with GRBM and yet contains certain amount of
information. To offset this problem, we propose a variant of GRBM called
multivariate Gaussian restricted Boltzmann machine.
### 2.3 Multivariate Gaussian Restricted Boltzmann Machine
Figure 1: Graphical model of a multivariate Gaussian restricted Boltzmann
machine. Links between sub-units in visible layer and hidden units are fully
connected. The interaction of the $k$th sub-unit in the $i$th visible unit
with $j$th hidden unit is modeled with $W_{i}^{kj}$.
The multivariate Gaussian restricted Boltzmann machine (MGRBM) is a natural
generalization of GRBM. The graphical model of a MGRBM is illustrated in
Figure 1. Compared with GRBM, in which each unit in visible layer is modeled
with a Gaussian distribution given the hidden layer, a MGRBM assumes that each
unit in visible layer is modeled with a multivariate Gaussian distribution
given the hidden layer. Similar to what $\sigma_{i}$ models in a GRBM, we
denote the covariance matrix of each unit in visible layer of a MGRBM with
$\Sigma_{i}$. Consider only the non-degenerate case, $\Sigma_{i}$ is a
positive-definite matrix and thus Cholesky decomposition can be applied:
$\Sigma_{i}=AA^{T}$. Since matrix $A$ is full-rank, we can denote $B=A^{-1}$.
With these notation, we can define the energy function of an MGRBM as:
$\displaystyle\begin{split}E(v,h;\theta)=\frac{1}{2}\sum_{i}(v_{i}-\mu_{i})^{T}B_{i}B_{i}^{T}(v_{i}-\mu_{i})\\\
-\sum_{i}v_{i}^{T}B_{i}W_{i}h-b^{T}h\end{split}$ (15)
Suppose the number of units in visible layer and hidden layer is $N_{v}$ and
$N_{h}$ respectively, and each unit in visible layer has $d$ dimension. Then
$v_{i},\mu_{i}$ each is a $d\times 1$ vector; $B_{i}$ each is a $d\times d$
matrix; $W_{i}$ each is a $d\times N_{h}$ matrix; $b$ and $h$ are both
$N_{h}\times 1$ vectors.
Similar to GRBM, we can also prove that:
$\displaystyle
p(v_{i}|h)=\mathcal{N}(\mu_{i}+(B_{i}^{T})^{-1}W_{i}h,(B_{i}^{T})^{-1}B_{i}^{-1})$
(16) $\displaystyle p(h_{j}=1|v)=sigmoid(b_{j}+\sum_{i}W_{ij}^{T}B_{i}v_{i})$
(17)
and learning algorithm is:
$\displaystyle\begin{split}\Delta\mu_{i}=\epsilon(<B_{i}B_{i}^{T}(v_{i}-\mu_{i})>_{data}\\\
-<B_{i}B_{i}^{T}(v_{i}-\mu_{i})>_{model})\end{split}$ (18)
$\displaystyle\Delta b=\epsilon(<h>_{data}-<h>_{model})$ (19)
$\displaystyle\Delta
W_{i}=\epsilon(<B_{i}^{T}v_{i}h^{T}>_{data}-<B_{i}^{T}v_{i}h^{T}>_{model})$
(20) $\displaystyle\begin{split}\Delta
B_{i}=\epsilon(<(v_{i}-\mu_{i})(v_{i}-\mu_{i})^{T}B_{i}-v_{i}h^{T}W_{i}^{T}>_{data}\\\
-<(v_{i}-\mu_{i})(v_{i}-\mu_{i})^{T}B_{i}-v_{i}h^{T}W_{i}^{T}>_{model})\end{split}$
(21)
In our experiment, we use persistent contrastive divergence (PCD) [9] to train
MGRBM. The algorithm can be written as:
$\displaystyle\begin{split}\Delta\mu_{i}=\epsilon(<B_{i}B_{i}^{T}(v_{i}-\mu_{i})>_{data}\\\
-<B_{i}B_{i}^{T}(v_{i}-\mu_{i})>_{fanta})\end{split}$ (22)
$\displaystyle\Delta b=\epsilon(<h>_{data}-<h>_{fanta})$ (23)
$\displaystyle\Delta
W_{i}=\epsilon(<B_{i}^{T}v_{i}h^{T}>_{data}-<B_{i}^{T}v_{i}h^{T}>_{fanta})$
(24) $\displaystyle\begin{split}\Delta
B_{i}=\epsilon(<(v_{i}-\mu_{i})(v_{i}-\mu_{i})^{T}B_{i}-v_{i}h^{T}W_{i}^{T}>_{data}\\\
-<(v_{i}-\mu_{i})(v_{i}-\mu_{i})^{T}B_{i}-v_{i}h^{T}W_{i}^{T}>_{fanta})\end{split}$
(25)
with $<\cdot>_{fanta}$ denotes average with respect to fantasy particles.
Notice that the problem for updating variances in a GRBM which we described in
section 2.2 still exists in MGRBM. In section 3.1 we will explain how we
address this problem in our experiments.
MGRBM is specially designed for speech data to address the problem of GRBM
described above. Typically, for the task of speech recognition with hybrid
HMM-neural network (NN) method, a context of several frames of acoustic
feature is used for each training case. Unlike GRBM, MGRBM can explicitly
model the evolving characteristics in each context. How a MGRBM can be used
for robust feature extraction is illustrated in figure 2. Concretely, suppose
the original acoustic feature for each frame has $D$ dimensions and each
context is chosen to be $C$ frames. Then the visible layer of MGRBM has $D$
units, each has dimension $C$; the $d$th dimension in the $c$th frame acoustic
feature corresponds to the $c$th dimension of $d$th unit.
Figure 2: Extract new feature from acoustic feature with multivariate Gaussian
restricted Boltzmann machine.
There are two reasons for the above setting. First, the correlation modeled
across each dimension of acoustic feature would act as a strong regularization
in temporal perspective if training data is much different from testing data.
Second, as is the case with GRBM, MGRBM also has conditional independence
assumption. Fortunately, this assumption is indeed satisfied if we use MFCC as
acoustic feature, for the step of discrete cosine transform already has the
effect of decorrelation.
### 2.4 Feature Learning for Robust Speech Recognition
Despite of its prevalence and huge success in phoneme recognition [10][11] and
LVCSR [1], deep neural network (DNN) stand alone is rarely used as a acoustic
model for robust speech recognition. We believe that this is due to the fact
that neural network is a discriminative model, whose optimization objective is
better discriminative power and lower classification error. However, for the
task of robust speech recognition, especially in mismatched scenario where the
training data is clean and the testing data is noisy, the power of such model
is greatly degraded due to its poor generalization over highly distorted data.
We assert that DNN performs significantly poorly in very noisy conditions than
GMM and we will prove it in our experiments later.
Generative models focus on modeling how the data is “generated”. Natural data
such as speech and image are usually high-dimensional, but all that make sense
occupy only a subspace (or lower-dimensional manifold). The learning process
of a generative model is essentially to find out this manifold by tuning all
its parameters. Background noise from real-life environment, we believe, is
substantially different from white noise, because the former kind of signal
contains certain characteristic shared with sounds that spread through the
air. For this reason, we consider it more appropriate to model speech with
generative models. So, we would like to find a model that can leverage such
advantage of generative models and yet escape from a model with strong
assumption like GMM, and RBM is indeed a such model.
As a generative model, RBM (and its variants) makes little assumption of data.
What is more important, it belongs to a family called product of experts (as
opposed to mixture of experts which GMM belongs to) [12]. This makes it
exponentially more powerful and less prone to over-fitting. Other than
modeling the distribution of data, it provides a natural way to transform
feature by applying (7)(14)(17). Since such transformation takes the whole
distribution of data into consideration, the learned feature tends to be much
more abstract and expressive.
## 3 Experiments
In this paper, we used the Aurora2 data set [13] for our experiments. In all
the experiments described below, we used only clean data set as training data
and whole test set as testing data.
We intentionally use acoustic models that are simple and comparable. Two
different kinds of acoustic model is utilized in our experiments : HMM-GMM and
hybrid HMM-DNN. With each kind of acoustic model, we compare the word error
rate (WER) of three kinds of feature : MFCC (12 coefficients + energy + delta
+ acceleration, 39-dimension), GRBM-extracted feature (G-feature) and MGRBM-
extracted feature (M-feature). We used standard Aurora2 setup described in
[13] for baseline system (GMM + MFCC). We use DNN simply because it is a good
classifier and it is much more natural to employ RBM-extracted feature.
### 3.1 RBM Training
We used two kinds of RBM in our experiment, one for feature extraction and the
other for pre-training of DNN.
GRBM and MGRBM are both trained to extract feature for comparison. For these
GRBMs and MGRBMs, PCD with stochastic gradient descent (SGD) was used. The
size of a mini-batch is 128 and no momentum was applied. The number of fantasy
particles was the same with the size of a mini-batch, and one full Gibbs
update was performed for each gradient estimate. For all the weights and
biases, learning rate was 0.001. For the updating of $B_{i}$ in (15) of an
MGRBM, learning rate was 0.0001. To avoid the problem of updating variances
described in section 2.2, we divide each $B_{i}$ by its trace after each
updating. This step makes all diagonal elements of $B_{i}$ fixed with one and
thus make the learning stable. All GRBM and MGRBM were trained for 400 epochs
and each with 1024 hidden units, which made G-feature and M-feature both has
1024 dimensions. The visible layer corresponds to a context of 9 frame of MFCC
feature, which makes the number of visible units of GRBM 351 and MGRBM
39$\times$9.
For all the GRBMs and RBMs that were used for pre-training DNNs, CD algorithm
(CD-1), SGD with batch size 128 and a momentum of 0.9 were used. We trained
all the RBMs with learning rate of 0.01 for 50 epochs and all the GRBMs with
learning rate of 0.001 for 100 epochs.
### 3.2 DNN Training
All DNNs in our experiments have 4 hidden layers, each hidden layer with 1024
units. DNNs were pre-trained with stacked RBM (section 3.1) as described in
[14] and fine-tuned with back-propagation algorithm with SGD as described in
[11]. The learning rate for back-propagation started from 1.0 and was halved
if an increase of substitution error on development set was observed during
the end of each epoch and all the weights roll-back to the end of last epoch.
When using MFCC feature, the input layer of DNN corresponds to 9 frame of
MFCC, which is 351 units. The output layer has 180 units, with each unit
corresponding to each state in HMM. All data to be trained with DNN are
normalized to have mean 0 and variance 1 with respect to each dimension.
### 3.3 Results
All results from our experiments are shown in Table 1. We averaged WER in all
test set across different noises. The first thing we should notice is that in
the lowest SNR condition, HMM-DNN performs consistently poorer than HMM-GMM,
which confirms our assertion in section 2.4. Those WER that exceeded 100% was
due to many substitution errors. From the last three columns, we can clearly
see the advantage of G-feature over MFCC and M-feature over G-feature.
Although the improvements does not seems to be obvious, the trend that the gap
between M-feature and G-feature increases with the decrease of SNR is still
distinguishable. The reason why the improvements of M-feature over G-feature
is marginal, we believe, is that the model which we trained is still not good
enough. All the learning rates, number of epochs and number of hidden units
are chosen heuristically, and 9 frame of context might be not sufficiently
long. So there is no reason to believe we have exploited the full potential of
MGRBM.
SNR | GMM | DNN
---|---|---
MFCC | G-feat (PCA) | M-feat (PCA) | MFCC | G-feat | M-feat
Clean | 1.05 | 1.94 | 2.75 | 1.27 | 0.76 | 0.82
20dB | 6.10 | 4.28 | 6.21 | 3.76 | 2.93 | 2.49
15dB | 15.18 | 8.50 | 13.67 | 9.09 | 6.47 | 5.64
10dB | 34.82 | 21.64 | 32.06 | 25.85 | 17.37 | 15.79
5dB | 61.27 | 49.80 | 59.75 | 58.05 | 41.87 | 40.99
0dB | 82.56 | 75.15 | 79.89 | 93.16 | 78.32 | 75.28
-5dB | 91.20 | 87.59 | 90.07 | 108.79 | 106.20 | 94.68
Table 1: Comparison of different features, all numbers are percentage of WER.
G-feat and M-feat here represents G-feature and M-feature respectively.
With HMM-GMM as acoustic model, the comparison is not so straightforward,
because training a 1024-dimensional GMM would leads to severe over-fitting. So
we reduced the dimension of G-feature and M-feature to 39 dimensions with
principal component analysis (PCA). Notice that this is actually not the right
thing to do, because the sum of top 39 eigenvalues only consists 91.8% and
69.7% of sum of all eigenvalues for G-feature and M-feature respectively.
Despite of such great loss of information, both G-feature and M-feature
performs better than MFCC for all SNR levels except the clean speech.
## 4 Conclusion and Discussion
In this paper, we briefly reviewed the definition and learning algorithm of
RBM and GRBM. We then propose a new variant of RBM called MGRBM by which we
would like to model the covariance of adjacent frames within each context.
After that we offered an explanation of why a feature learned with RBM (and
its variants) might be able to enhance the performance of robust speech
recognition over original feature. Finally we performed our experiments on
Aurora2 and showed that feature that learned with GRBM and MGRBM would indeed
improve the average accuracy across environment in every SNR condition over
original MFCC feature.
Throughout the process of feature learning with GRBM, virtually nothing is
presupposed. This makes it adaptable with any feature generated from a front-
end. Aside from training a GRBM, which can be done off-line, the extra cost of
the feature transformation is merely a multiplication with a matrix. From this
perspective, feature learning with GRBM is similar to the TANDEM system [15],
but the latter is purely supervised. Hence many advantages can be gained if
lots of data is accessible but little is labeled. What is more important,
GRBM-extracted feature can be used for the TANDEM System seamlessly, which
makes TANDEM system semi-supervised and thus further enhancement of
performance can be achieved.
## 5 Acknowledgements
This work is partially supported by the National Basic Research Program (973
Program) of China(2012CB316401), the National Natural Science Foundation of
China (60928005, 60805008, 60931160443 and 61003094), the Ph.D. Programs
Foundation of Ministry of Education of China (200800031015), the Upgrading
Plan Project of Shenzhen Key Laboratory and the Science and Technology R&D
Funding of the Shenzhen Municipal.
## References
* [1] G. Dahl, D. Yu, L. Deng, , and A. Acero, “Large vocabulary continuous speech recognition with context-dependent dbn-hmms,” in _Proc. ICASSP_ , 2011.
* [2] G. E. Hinton, L. Deng, D. Yu, G. Dahl, A. Mohamed, N. Jaitly, A. Senior, V. Vanhoucke, P. Nguyen, T. Sainathand, and B. Kingsbury, “Deep neural networks for acoustic modeling in speech recognition,” _IEEE Signal Processing Magazine_ , 2012.
* [3] ETSI, “Advanced front-end feature extraction algorithm,” ETSI ES 202 050, Tech. Rep., 2007.
* [4] Y. Bengio, A. Courville, and P. Vincent, “Representation learning: A review and new perspectives,” arXiv:1206.5538 [cs.LG].
* [5] P. Smolensky, “Information processing in dynamical systems: Foundations of harmony theory,” _Parallel Distributed Processing: Explorations in the Microstructure of Cognition_ , vol. 1, 1986.
* [6] G. E. Hinton, “Training products of experts by minimizingc ontrastive divergence,” _Neural Computation_ , vol. 14, pp. 1771–1800, 2002.
* [7] G. E. Hinton and R. Salakhutdinov, “Reducing the dimensionality of data with neural networks,” _Science_ , vol. 313, pp. 504–507, 2006.
* [8] G. W. Taylor, G. E. Hinton, and S. T. Roweis, “Two distributed-state models for generating high-dimensional time series,” _The Journal of Machine Learning Research_ , vol. 12, pp. 1025–1068, 2011.
* [9] T. Tieleman, “Training restricted boltzmann machines using approximations to the likelihood gradient,” in _Proc. of ICML_ , 2008.
* [10] A. Mohamed, G. E. Dahl, and G. E. Hinton, “Deep belief networks for phone recognition,” in _NIPS Workshop on Deep Learning for Speech Recognition and Related Applications_ , 2009.
* [11] A. Mohamed, G. Dahl, and G. Hinton, “Acoustic modeling using deep belief networks,” _IEEE Trans. on Audio, Speech and Language Processing_ , 2011\.
* [12] G. E. Hinton, “Products of experts,” _ICANN_ , vol. 1, pp. 1–6, 1999.
* [13] H. G. Hirsch and D. Pierce, “The aurora experimental framework for the performance evaluation of speech recognition systems under noise conditions,” in _Proc. ISCA ITRW ASR2000_ , 2000.
* [14] G. E. Hinton, S. Osindero, and Y. Teh, “A fast learning algorithm for deep belief nets,” _Neural Computation_ , vol. 18, pp. 1527–1554, 2006.
* [15] H. Hermansky, D. P. W. Ellis, and S. Sharma, “Tandem connectionist feature extraction for conventional hmm systems,” in _Proc. ICASSP_ , 2000, pp. 1635–1638.
|
arxiv-papers
| 2013-09-23T13:51:28 |
2024-09-04T02:49:51.446285
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Xin Zheng, Zhiyong Wu, Helen Meng, Weifeng Li, Lianhong Cai",
"submitter": "Xin Zheng",
"url": "https://arxiv.org/abs/1309.6176"
}
|
1309.6266
|
# ENERGY OF SIGNED DIGRAPHS
S. Pirzada and Mushtaq A. Bhat Department of Mathematics
University of Kashmir
Srinagar, Hazratbal 190006
India [email protected]; [email protected]
[email protected]
###### Abstract.
In this paper we extend the concept of energy to signed digraphs. We obtain
Coulson’s integral formula for energy of signed digraphs. Formulae for
energies of signed directed cycles are computed and it is shown that energy of
non cycle balanced signed directed cycles increases monotonically with respect
to number of vertices. Characterization of signed digraphs having energy equal
to zero is given. We extend the concept of non complete extended $p$ sum (or
briefly, NEPS) to signed digraphs. An infinite family of equienergetic signed
digraphs is constructed. Moreover, we extend McClelland’s inequality to signed
digraphs and also obtain sharp upper bound for energy of signed digraph in
terms the number of arcs. Some open problems are also given at the end.
###### Key words and phrases:
Energy of signed digraph, NEPS,McClelland’s inequality.
###### 2000 Mathematics Subject Classification:
fill it,
###### 2010 Mathematics Subject Classification:
05C50, 05C22,05C76.
## 1\. introduction
A signed digraph (or briefly sidigraph) is defined to be a pair $S=(D,\sigma)$
where $D=(V,\mathscr{A})$ is the underlying digraph and
$\sigma:\mathscr{A}\rightarrow\\{-1,1\\}$ is the signing function. The sets of
positive and negative arcs of $S$ are respectively denoted by
$\mathscr{A}^{+}(S)$ and $\mathscr{A}^{-}(S)$. Thus
$\mathscr{A}(S)=\mathscr{A}^{+}(S)\cup\mathscr{A}^{-}(S)$. A sidigraph is said
to be homogeneous if all of its arcs have either positive sign or negative
sign, otherwise heterogeneous. A sidigraph without orientations of its arcs
(when it is regarded as a simple undirected sigraph) is known as its
underlying sigraph and is denoted by $S^{u}$.
Two vertices are adjacent if they are connected by an arc. A path of length
$n-1$ $(n\geq 2)$, denoted by $P_{n}$, is a sidigraph on $n$ vertices
$v_{1},v_{2},\cdots,v_{n}$ with $n-1$ signed arcs $(v_{i},v_{i}+1)$. A cycle
of length $n$, is a sidigraph having vertices $v_{1},v_{2},\cdots,v_{n}$ and
signed arcs $(v_{i},v_{i}+1)$ $i=1,2,\cdots,n-1$ and $(v_{n},v_{1})$. A
sidigraph is linear if each vertex has both indegree and outdegree equal to
one. The sign of a sidigraph is defined as the product of signs of its arcs. A
sidigraph is said to be positive (negative) if its sign is positive (negative)
i.e., it contains an even (odd) number of negative arcs. A signed digraph is
said to be all-positive (respectively, all negative) if all its arcs are
positive (negative). A signed digraph is said to be cycle balanced if each of
its cycles is positive, otherwise non cycle balanced. Throughout this paper we
call cycle balanced cycle a positive cycle and non cycle balanced cycle a
negative cycle and respectively denote them by $C_{n}$ and ${\bf C}_{n}$,
where $n$ is number of vertices.
A sidigraph is symmetric if
$(u,v)\in\mathscr{A}^{+}(S)~{}~{}(or~{}~{}\mathscr{A}^{-}(S))$ then
$(v,u)\in\mathscr{A}^{+}(S)~{}~{}(or~{}~{}\mathscr{A}^{-}(S))$, where $u,v\in
V(S)$. A one to one correspondence between sigraphs and symmetric sidigraphs
is given by $S\rightsquigarrow\overleftrightarrow{S}$, where
$\overleftrightarrow{S}$ has the same vertex set as that of sigraph $S$, and
each signed edge $(u,v)$ is replaced by a pair of symmetric arcs $(u,v)$ and
$(v,u)$ both with same sign as that of edge $(u,v)$. Under this correspondence
a sigraph can be identified with a symmetric sidigraph. A sidigraph is said to
be skew symmetric if its adjacency matrix is skew symmetric. A skew symmetric
sidigraph is obtained by replacing each edge of sigraph by a pair of opposite
oriented and opposite signed arcs. We denote a skew symmetic digraph of order
$n$ by ${\bf S}_{n}$.
Two sidigraphs (sigraphs) on same number of vertices are said to be cospectral
if they have same spectrum. Acharya, Gill, Patwardhan [2] introduced the
concept of Quasicospectrality in digraphs (graphs). Two digraphs (graphs) $G$
and $H$ of same order are said to be quasicospectral if there exist sidigraphs
(sigraphs) $S$ and $S^{\prime}$ on $G$ and $H$ respectively which are
cospectral. The weighted directed graph $S$ of an $n\times n$ matrix
$M=(m_{ij})$ of reals consists of $n$ vertices, with vertex $i$ joined by a
directed arc with weight $m_{ij}$ to vertex $j$ if and only if $m_{ij}$ is
non-zero. In case the matrix consists of entries $-1,0$ and $1$, then we get a
sidigraph. Thus there is a one to one correspondence between the set of
integral $(-1,0,1)$-matrices of order $n$ and set of sidigraphs of order $n$.
The cartesian product (or sum) of two sidigraphs
$S_{1}=(V_{1},A_{1},\sigma_{1})$ and $S_{2}=(V_{2},A_{2},\sigma_{2})$ denoted
by $S_{1}\times S_{2}$ is the sidigraph $(V_{1}\times V_{2},A,\sigma)$, where
the arc set is that of the cartesian product of unsigned digraphs and the sign
function is defined by
$\sigma((u_{i},v_{j}),(u_{k},v_{l}))=\left\\{\begin{array}[]{lr}\sigma_{1}(u_{i},u_{k}),&\mbox{if
$j=l$},\\\ \sigma_{2}(v_{j},v_{l}),&\mbox{if $i=k.$}\end{array}\right.$
The Kronecker product (strong product or conjunction) of two sidigraphs
$S_{1}=(V_{1},A_{1},\sigma_{1})$ and $S_{2}=(V_{2},A_{2},\sigma_{2})$ denoted
by $S_{1}\otimes S_{2}$ is the sidigraph $(V_{1}\times V_{2},A,\sigma)$, where
arc set is the arc set of underlying unsigned digraphs and the sign function
is defined by
$\sigma((u_{i},v_{j}),(u_{k},v_{l}))=\sigma_{1}(u_{i},u_{k})\sigma_{2}(v_{j},v_{l}).$
The adjacency matrix of a sidigraph $S$ whose vertices are
${v_{1},v_{2},\cdots,v_{n}}$ is the $n\times n$ matrix $A(S)=(a_{ij})$, where
$a_{ij}=\left\\{\begin{array}[]{lr}\sigma(v_{i},v_{j}),&\mbox{if there is an
arc from $v_{i}$ to $v_{j},$}\\\ 0,&\mbox{otherwise.}\end{array}\right.$
The characteristic polynomial $|xI-A(S)|$ of the adjacency matrix $A(S)$ of
sidigraph $S$ is called the characteristic polynomial of $S$ and is denoted by
$\phi_{S}(x)$. The eigenvalues of $A(S)$ are called eigenvalues of $S$. The
following is the coefficient Theorem for sidigraphs [2]
Theorem 1.1. If $S$ be a sidigraph with characteristic polynomial
$\phi_{S}(x)=x^{n}+c_{1}x^{n-1}+\cdots+c_{n-1}x+c_{n}$
then
$c_{i}=\sum\limits_{L\in\pounds_{i}}(-1)^{p(L)}\prod\limits_{Z\in c(L)}s(Z),$
for all $i=1,2,\cdots,n$, where $\pounds_{i}$ is the set of all linear
subdigraphs $L$ of $S$ of order $i$, $p(L)$ denotes number of components of
$L$ and $c(L)$ denotes the set of all cycles of $L$ and $s(Z)$ the sign of
cycle $S$.
If $S$ is an undirected sigraph then $S$ can be viewed as a digraph $\bar{S}$
by identifying each signed edge of $S$ with a directed cycle of length $2$
with each arc having same sign as that of the corresponding edge. For
undirected sigraphs Theorem $1.1$ takes the following form [10].
Theorem 1.2. If $S$ is a sigraph with characteristic polynomial
$\phi_{S}(x)=x^{n}+c_{1}x^{n-1}+\cdots+c_{n-1}x+c_{n}$
then
$c_{i}=\sum\limits_{L\in\pounds_{i}}(-1)^{p(L)}2^{c(L)}\prod\limits_{Z\in
c(L)}s(Z),$
for all $i=1,2,\cdots,n$, where $\pounds_{i}$ is the set of all basic figures
$L$ of $S$ of order $i$, $p(L)$ denotes number of components of $L$ and $c(L)$
the number of cyclic components of $L$.
The spectral criterion for cycle balance of sidigraphs (sigraphs) given by
Acharya is as follows.
Theorem 1.3[1, 2]. A sidigraph (sigraph) is cycle balanced (balanced) if and
only if it is cospectral with the underlying unsigned digraph (graph).
## 2\. Energy of sidigraphs
For spectra and energy of graphs and digraphs see [5, 6, 11, 12, 15]. Germina,
Hameed and Zaslavsky [7] defined energy of a sigraph to be the sum of absolute
values of sigraph eigenvalues. In this section, we extend the concept of
energy to sidigraphs. Unlike sigraphs the adjacency matrix of a sidigraph need
not be real symmetric, so eigenvalues can be complex numbers.
Definition 2.1. Let $S$ be a sidigraph of order $n$ having eigenvalues
$z_{1},z_{2},\cdots,z_{n}$. The energy of $S$ is defined as
$E(S)=\sum\limits_{j=1}^{n}|\Re z_{j}|,$
where $\Re z_{j}$ denotes real part of complex number $z_{j}.$
If $S$ is a signed graph and $\overleftrightarrow{S}$ be its symmetric
sidigraph, then clearly $A(S)=A(\overleftrightarrow{S})$ and so
$E(S)=E(\overleftrightarrow{S})$. In this way, definition $2.1$ generalizes
the concept of energy of undirected sigraphs.
Example 2.2. Let $S$ be a sidigraph shown in Figure $1$. Dotted arrows denote
negative arcs and plane ones denote positive arcs. Clearly, $S$ is non cycle
balanced sidigraph. By Theorem $1.1$, the characteristic polynomial of $S$ is
$\phi_{S}(x)=x^{10}+x^{7}=x^{7}(x^{3}+1)$. The eigenvalues of $S$ are
$0^{7},1,\frac{1\pm\sqrt{3}\iota}{2}$, where $\iota=\sqrt{-1}$, so $E(S)=2$.
SFigure 1
Example 2.3. Let $S$ be an acyclic sidigraph. Then by Theorem $1.1$, the
characteristic polynomial of $S$ is $\phi(S)=x^{n}$, so that
$spec(S)=\\{0^{n}\\}$ and hence $E(S)=0$.
Example 2.4. Consider ${\bf S}_{n}$, the skew symmetric sidigraph on $n\geq 2$
vertices, then by Theorem $1.1$, $\phi_{{\bf S}_{n}}(x)=x^{n-2}(x^{2}+(n-1))$
so that eigenvalues are $0^{n-2},\pm\iota\sqrt{n-1}$. Thus all skew symmmetric
sidigraphs are cospectral with energy equal to zero.
Example 2.5. If $S$ is the sidirected cycle on $n$ vertices, then the
characteristic polynomial of $S$ is $\phi(S)=x^{n}+(-1)^{[s]}$, where the
symbol $[s]$ is defined as $[s]=1$ or $0$ according as $S$ is positive or
negative. If $S=C_{n}$, then eigenvalues are $e^{\frac{2\iota
j\pi}{n}},~{}j=0,1,\cdots,n-1$. So,
$E(S)=\sum\limits_{j=0}^{n-1}|\cos(\frac{2j\pi}{n})|$. If $S={\bf C}_{n}$,
then the eigenvalues are $e^{\frac{\iota(2j+1)\pi}{n}},~{}j=0,1,\cdots,n-1$.
Therefore, $E(S)=\sum\limits_{j=0}^{n-1}|\cos(\frac{(2j+1)\pi}{n})|$. In
particular if $S={\bf C}_{4}$, the eigenvalues are
$\frac{1\pm\iota}{\sqrt{2}},\frac{-1\pm\iota}{\sqrt{2}}$ and $E(S)=2\sqrt{2}$.
Example 2.6. Let $S$ be a sidigraph having $n$ vertices and unique cycle of
length $r$, where $2\leq r\leq n$. Then by Theorem $1.1$,
$\phi_{S}(x)=x^{n}+(-1)^{[s]}x^{r}=x^{r}(x^{n-r}+(-1)^{[s]})$, where the
symbol $[s]$ is defined as $[s]=1$ or $0$ according as $S$ is cycle balanced
or non cycle balanced. Clearly, energy equals to the energy of the unique
cycle.
Given $t$ sidigraphs $S_{1},S_{2},\cdots S_{t}$, their direct product denoted
by $S_{1}\oplus S_{2}\oplus\cdots\oplus S_{t}$ is the sidigraph with
$V_{S_{1}\oplus S_{2}\oplus\cdots\oplus
S_{t}}=\bigcup\limits_{j=1}^{t}V(S_{j})$ and arc set $\mathscr{A}_{S_{1}\oplus
S_{2}\oplus\cdots\oplus S_{t}}=\bigcup\limits_{j=1}^{t}\mathscr{A}(S_{j})$.
Theorem 2.7. Let $S$ be a sidigraph on n vertices and
$S_{1},S_{2},\cdots,S_{k}$ be its strong components. Then
$E(S)=\sum\limits_{j=1}^{k}E(S_{j}).$
Proof. Let $Y=\\{a\in A(S):a\notin c(S)\\}$, where $c(S)$ is the set of all
cycles of $S$. By Theorem $1.1$, $\phi_{S}(x)=\phi_{S-Y}(x)$, where $S-Y$ is
the sidigraph obtained from $S$ by deleting the non-cyclic arcs. Clearly,
$S-Y=S_{1}\oplus S_{2}\oplus\cdots\oplus S_{k}$ and adjacency matrix of
sidigraph $S-Y$ is in block diagonal form with diagonal blocks as the
adjacency matrices of strong components. Therefore
$\phi_{S-Y}(x)=\phi_{S_{1}}(x)\phi_{S_{2}}(x)\cdots\phi_{S_{k}}(x)$ and so
$E(S)=\sum\limits_{j=1}^{k}E(S_{j})$. ∎
Remark 2.8. From Theorem $1.1$,
$c_{i}=\sum\limits_{L\in\pounds_{i}}(-1)^{p(L)},i=1,2,\cdots,n$. Clearly, this
sum contains positive and negative ones.
$+1$ arises if and only if
$(a)$ Number of components of $L\in\pounds_{i}$ are odd and $s(L)<0$. We call
such linear sidigraphs as type $a$ linear sidigraphs.
$(b)$ Number of components of $L\in\pounds_{i}$ is even and $s(L)>0$. We call
such linear sidigraphs as type $b$.
$-1$ will occur if and only if
$(c)$ Number of components of $L\in\pounds_{i}$ is odd and $s(L)>0$. We call
such linear sidigraphs as type $c$.
$(d)$ Number of components of $L\in\pounds_{i}$ is even and $s(L)<0$. We call
such linear sidigraphs as type $d$.
From above we observe that $c_{i}=0$ if and only if either $S$ is acyclic or
for each $i$, number of linear sidigraphs of order $i$ of type $a$ or type $d$
is equal to the number of linear sidigraphs of order $i$ of type $b$ or type
$c$.
An immediate consequence of Remark $2.7$ is the following Lemma.
Lemma 2.9. An integral $(-1,0,1)$ matrix is nilpotent if and only if its
underlying sidigraph $S$ is either acyclic or in $S$, for each
$i=1,2,\cdots,n$, the number of linear sidigraphs of order $i$ of type $a$ or
type $d$ is equal to number of linear sidigraphs of order $i$ of type $b$ or
type $c$. ∎
Unlike unsigned strong component, energy of a signed directed strong component
can be zero, for example, sidigraph $S_{1}$ in Figure 2. The following result
characterizes sidigraphs having energy equal to zero.
Theorem 2.10. Let $S$ be a sidigraph of order $n$. Then $E(S)=0$ if and only
if $S$ satisfies one of the following conditions $(I)$ $S$ is acyclic $(II)$
$S$ is skew symmetric $(III)$ for each $i=1,2,\cdots,n$, the number of linear
sidigraphs of order $i$ of type $a$ or type $d$ is equal to number of linear
sidigraphs of order $i$ of type $b$ or type $c$.
Proof. Let $S$ be a sidigraph of order $n$. If $S$ is acyclic, then by Lemma
$2.9$, $\phi_{S}(x)=x^{n}$ and so $E(S)=0.$ If $S$ is skew symmetric
sidigraph, then the eigenvalues of $S$ are $0^{n-2},\pm\iota\sqrt{n-1}$,
therefore $E(S)=0$. If $S$ satisfies $(III)$, then by Lemma $2.9$,
$\phi_{S}(x)=x^{n}$, so that $E(S)=0.$
Conversely, if $E(S)=0$, then either eigenvalues of $S$ are
$0^{n-2},\pm\iota\sqrt{n-1}$ or all equal to zero. Using the fact that a real
matrix is skew symmetric if and only if all its eigenvalues are of the form
$\pm\iota\alpha$, where $\alpha\in\mathbb{R}$ and by Lemma $2.9$, the result
follows. ∎
## 3\. Computation of energy of signed directed cycles
We first calculate energy formulae for positive cycles. Let $C_{n}$ be a
positive cycle on $n\geq 2$ vertices. The characteristic polynomial of $C_{n}$
is $\phi_{C_{n}}(x)=x^{n}-1$, so that the eigenvalues are $e^{\frac{2\pi\iota
j}{n}},j=0,1,\cdots,n-1,$ $\iota=\sqrt{-1}$. Consequently energy of $C_{n}$ is
$E(C_{n})=\sum\limits_{j=0}^{n-1}|\cos\frac{2j\pi}{n}|.$
Given a positive integer $n$, it has one of the forms $4k$, or $2k+1$, or
$4k+2$, where $k\geq 0$.
If $n=4k$, then
$\displaystyle E(C_{n})$
$\displaystyle=\sum\limits_{j=0}^{4k-1}|\cos\frac{2j\pi}{4k}|=\sum\limits_{j=0}^{4k-1}|\cos\frac{j\pi}{2k}|=2\sum\limits_{j=0}^{2k-1}|\cos\frac{j\pi}{2k}|$
$\displaystyle=2+4\sum\limits_{j=1}^{k-1}\cos\frac{j\pi}{2k}=2+4\\{\frac{-1}{2}+\frac{\sin\frac{(k-\frac{1}{2})\pi}{2k}}{2\sin\frac{\pi}{4k}}\\}=2\cot\frac{\pi}{n}.$
If $n=2k+1$, then
$\displaystyle E(C_{n})$
$\displaystyle=\sum\limits_{j=0}^{2k}|\cos\frac{2j\pi}{2k+1}|=1+2\sum\limits_{j=1}^{k}|\cos\frac{2j\pi}{2k+1}|=1+2\sum\limits_{j=1}^{k}\cos\frac{j\pi}{2k+1}$
$\displaystyle=1+2\\{\frac{-1}{2}+\frac{\sin\frac{(k+\frac{1}{2})\pi}{2k+1}}{2\sin\frac{\pi}{2(2k+1)}}\\}=2\csc\frac{\pi}{2n}.$
If $n=4k+2$ then
$\displaystyle E(C_{n})$
$\displaystyle=\sum\limits_{j=0}^{4k+1}|\cos\frac{2j\pi}{4k+2}|=\sum\limits_{j=0}^{4k+1}|\cos\frac{j\pi}{2k+1}|=2\sum\limits_{j=0}^{2k}|\cos\frac{j\pi}{2k+1}|$
$\displaystyle=2+4\sum\limits_{j=1}^{k}\cos\frac{j\pi}{2k+1}=2+4\\{\frac{-1}{2}+\frac{\sin\frac{(k+\frac{1}{2})\pi}{2k+1}}{2\sin\frac{\pi}{4k+2}}\\}=2\csc\frac{\pi}{n}.$
We now calculate exact formulae for the energy of negative cycles of length
$n$. Let ${\bf C}_{n}$ denote the negative cycle with $n$ vertices. Then
$\phi_{{\bf C}_{n}}(x)=x^{n}+1$ and so $Spec({\bf
C}_{n})=e^{\frac{(2j+1)\pi\iota}{n}},j=0,1,\cdots,n-1$, $\iota=\sqrt{-1}$.
Therefore energy is given by
$E({\bf C}_{n})=\sum\limits_{j=0}^{n-1}|\cos\frac{(2j+1)\pi}{n}|.$
If $n=4k$, then
$\displaystyle E({\bf C}_{n})$
$\displaystyle=\sum\limits_{j=0}^{4k-1}|\cos\frac{(2j+1)\pi}{4k}|=2\sum\limits_{j=0}^{2k-1}|\cos\frac{(2j+1)\pi}{4k}|=4\sum\limits_{j=0}^{k-1}\cos\frac{(2j+1)\pi}{4k}$
$\displaystyle=4\\{\cos\frac{\pi}{4k}+\cos\frac{3\pi}{4k}+\cdots+\cos\frac{(2k-1)\pi}{4k}\\}=4\frac{\cos(\frac{\pi}{4k}+\frac{k-1}{2}\frac{2\pi}{4k})\sin
k\frac{2\pi}{8k}}{\sin\frac{2\pi}{8k}}=2\csc\frac{\pi}{n}.$
If $n=4k+2$, then
$\displaystyle E({\bf C}_{n})$
$\displaystyle=\sum\limits_{j=0}^{4k+1}|\cos\frac{(2j+1)\pi}{4k+2}|=4\sum\limits_{j=0}^{k-1}\cos\frac{(2j+1)\pi}{4k+2}$
$\displaystyle=4\\{\cos\frac{\pi}{4k+2}+\cos\frac{3\pi}{4k+2}+\cdots+\cos\frac{(2k-1)\pi}{4k+2}\\}=2\cot\frac{\pi}{n}.$
If $n=2k+1$, then since $-1$ is the eigenvalue of ${\bf C}_{n}$, we have
$spec({\bf C}_{n})$ $=-spec(C_{n})$, and so $E({\bf C}_{n})=E(C_{n})$.
Summarizing, all the above cases can be written as follows:
$E(C_{n})=\left\\{\begin{array}[]{lr}2\cot\frac{\pi}{n},&\mbox{if $n=4k$},\\\
\csc\frac{\pi}{2n},{}{}{}{}&\mbox{if $n=4k+1$~{}or~{}$n=4k+3$},\\\
2\csc\frac{\pi}{n},&\mbox{if $n=4k+2$}\\\ \end{array}\right.$
and
$E(\bf{C}_{n})=\left\\{\begin{array}[]{lr}2\csc\frac{\pi}{n},&\mbox{if
$n=4k$},\\\ \csc\frac{\pi}{2n},{}{}{}{}&\mbox{if $n=4k+1$~{}or $n=4k+3$},\\\
2\cot\frac{\pi}{n},&\mbox{if $n=4k+2$}.\\\ \end{array}\right.$
Pena and Rada [15] proved that the energy of directed unsigned cycles
increases monotonically with respect to order $n$. From energy formulae for
positive and negative sidirected cycles, the following two results are
immediate.
Theorem 3.1. Among all non cycle balanced unicyclic sidigraphs on $n$
vertices, the cycle has the largest energy. Moreover, energy of negative
cycles increases monotonically with respect to the order. Minimal energy is
attained in ${\bf C}_{2}$, a negative cycle of order $2$. ∎
Theorem 3.2. Energy of positive and negative cycles satisfy the following:
$(1)$ Energy of positive cycle of odd order equals energy of negative cycle of
same order.
$(2)$ Energy of negative cycle of even order is greater than energy of
positive cycle of same order if and only if $n=4k$.
$(3)$ Energy of negative cycle of even order is less than energy of positive
cycle of same order if and only if $n=4k+2$.
One of the fundamental results in theory of graph energy is Coulsons integral
formula. We now state Coulsons integral formula for energy of sidigraphs. The
proof is similar as in digraphs.
Theorem 3.3. Let $S$ be a sidigraph with $n$ vertices having characteristic
polynomial $\phi_{S}(x)$. Then
$E(S)=\sum\limits_{j=1}^{n}|\Re
z_{j}|=\frac{1}{\pi}\int\limits_{-\infty}^{\infty}(n-\frac{\iota
x\phi^{\prime}_{S}(\iota x)}{\phi_{S}(\iota x)})dx,$
where $z_{1},z_{2},\cdots,z_{n}$ are the eigenvalues of sidigraph $S$ and
$\int\limits_{-\infty}^{\infty}F(x)dx$ denotes principle value of the
respective integral.
The Coulson’s formula given above is a motivation to define energy of
sidigraphs to be the sum of absolute values of real parts of eigenvalues.
Example 3.4. Consider the cycle ${\bf C}_{4}$, the characteristic polynomial
is $\phi_{{\bf C}_{4}}(x)=x^{4}+1$ and hence
$E({\bf C}_{4})=\frac{1}{\pi}\int\limits_{-\infty}^{\infty}[4-\frac{4\iota
x(\iota x)^{3}}{(\iota x)^{4}+1}]dx$
$=\frac{1}{\pi}\int\limits_{-\infty}^{\infty}\frac{4}{x^{4}+1}dx$
$=\frac{4}{\pi}\frac{\pi}{2\sin\frac{\pi}{4}}=2\sqrt{2}$, as calculated in
example $2.5$.
An immediate consequence of Coulson’s integral formula is the following
observation, the proof of which is similar to the case in unsigned digraphs.
Corollary 3.5. If $S$ is a sidigraph on $n$ vertices, then
$E(S)=\frac{1}{\pi}\int\limits_{-\infty}^{\infty}\frac{1}{x^{2}}\log(x^{n}\phi_{G}(\frac{\iota}{x})dx.$
Pena and Rada [15] considered the problem of increasing property of the energy
in digraphs, It is natural to consider the same problem for sidigraphs.
Let ${S}_{n,h}$ denote the set of sidigraphs with $n$ vertices and every cycle
having length $h$. First we shall calculate characteristic polynomial of such
sidigraphs.
Theorem 3.6. If $S\in S_{n,h}$, then
$\phi_{S}(x)=x^{n}+\sum\limits_{j=1}^{\lfloor{\frac{n}{h}}\rfloor}(-1)^{k}c^{*}(S,kh)x^{n-kh},$
where $c^{*}(S,kh)=$ number of positive linear sidigraphs of order $kh$ $-$
number of negative linear sidigraphs of order $kh$, for every
$k=1,2,\cdots,{\lfloor{\frac{n}{h}}\rfloor}$.
Proof. By Theorem $1.1$, the coefficient
$c_{kh}=\sum\limits_{L\in\pounds_{kh}}(-1)^{p(L)}\prod\limits_{Z\in
c(L)}s(Z)$. As each cycle of $S$ has length $h$, it follows that
$c_{kh}=(-1)^{k}c^{*}(S,kh)$, where $c^{*}(S,kh)=$ number of positive linear
sidigraphs consisting of $k$ cycles each of length $h$ $-$ number of negative
linear sidigraphs consisting of $k$ cycles each of length $h$. Clearly,
$c_{j}=0$ if $j$ is not a multiple of $h$, because $\pounds_{j}=\phi$ in that
case. ∎
Definition 3.7. Let $S^{*}_{n,h}=\\{S\in S_{n,h}:c^{*}(S,kh)\geq
0,k=1,2,\cdots,\lfloor\frac{n}{h}\rfloor\\}$. Define a quasi order relation
over $S^{*}_{n,h}$ as follows. Let $S_{1}$ and $S_{2}$ be two elements of
$S^{*}_{n,h}$. Define $S_{1}\preceq S_{2}$ if for all
$k=1,2,\cdots,\lfloor\frac{n}{h}\rfloor$, $c^{*}(S_{1},kh)\leq
c^{*}(S_{2},kh)$. If $S_{1}\preceq S_{2}$ and there exists $k$ such that
$c^{*}(S_{1},kh)<c^{*}(S_{2},kh)$ then $S_{1}\prec S_{2}$. Clearly, this is a
reflexive and transitive relation over $S^{*}_{n,h}$.
Theorem 3.8. Let $h$ be an integer of the form $4t+2,~{}t\geq 0$. The energy
increases with respect to quasi-order relation defined over $S^{*}_{n,h}$,
i.e., if $S_{1},~{}S_{2}\in S^{*}_{n,h}$ then $S_{1}\prec S_{2}\Rightarrow
E(S_{1})<E(S_{2}).$
Proof. Let $H\in S^{*}_{n,h}$, by Theorem $3.6$, we have
$\phi_{H}(x)=x^{n}+\sum\limits_{j=1}^{\lfloor{\frac{n}{h}}\rfloor}(-1)^{k}c^{*}(H,kh)x^{n-kh},$
so that
$\phi_{H}(\frac{\iota}{x})=\frac{\iota^{n}}{x^{n}}+\sum\limits_{j=1}^{\lfloor{\frac{n}{h}}\rfloor}(-1)^{k}c^{*}(H,kh)\frac{\iota^{n-kh}}{x^{n-kh}}$
$=\frac{\iota^{n}}{x^{n}}(1+\sum\limits_{j=1}^{\lfloor{\frac{n}{h}}\rfloor}(-1)^{k}c^{*}(H,kh)x^{kh}\iota^{-k(4t+2)})$
$=\frac{\iota^{n}}{x^{n}}(1+\sum\limits_{j=1}^{\lfloor{\frac{n}{h}}\rfloor}c^{*}(H,kh)x^{kh}).$
Applying Corollary $3.5$, we have
$E(S)=\frac{1}{\pi}\int\limits_{-\infty}^{\infty}\frac{1}{x^{2}}\log(x^{n}\frac{\iota^{n}}{x^{n}}(\sum\limits_{j=1}^{\lfloor{\frac{n}{h}}\rfloor}c^{*}(H,kh)x^{kh}))dx=$
$\frac{1}{\pi}\int\limits_{-\infty}^{\infty}\frac{1}{x^{2}}\log(\iota^{n}(1+\sum\limits_{j=1}^{\lfloor{\frac{n}{h}}\rfloor}c^{*}(H,kh)x^{kh}))dx$
$=\frac{1}{\pi}\int\limits_{-\infty}^{\infty}\frac{1}{x^{2}}\log(1+\sum\limits_{j=1}^{\lfloor{\frac{n}{h}}\rfloor}c^{*}(H,kh)x^{kh})dx$,
as
$\frac{1}{\pi}p.v.\int\limits_{-\infty}^{+\infty}\log(\iota^{n})\frac{dx}{x^{2}}=0$,
where $p.v.$ stands for principal value of Cauchy’s integral. It is clear from
above energy expression that energy increases with respect to quasi-order
relation defined over $S^{*}_{n,h}$. ∎
## 4\. NEPS in sidigraphs
Recall that [13] Kronecker product of two matrices $A=(a_{ij})_{r\times s}$
and $B=(b_{ij})_{t}\times u$ denoted by $A\otimes B$ is matrix of order
$rt\times su$ obtained by replacing each entry $a_{ij}$ of $A$ by a block
$a_{ij}B$. Thus $A\otimes B$ consists of all $rtsu$ possible products of an
entry of $A$ with an entry of $B$. The Kronecker product is a componentwise
operation, i.e., $(A\otimes B)(C\otimes D)=(AC)\otimes(BD)$, provided the
products $AC$ and $BD$ exist. This operation is also associative, so we can
define the multiple product $A_{1}\otimes A_{2}\otimes\cdots\otimes A_{m}$.
Let order of $A_{i}$ be $r_{i}\times s_{i}$. We index elements of $A_{i}$ by
$a_{i};jk$, and those of multiple product by a pair of $m-$tuples, a row index
$j=(j_{1},j_{2},\cdot,j_{m})$ and a column index
$k=(k_{1},k_{2},\cdots,k_{m})$, where $1\leq j_{i}\leq r_{i}$ and $1\leq
k_{i}\leq s_{i}$. The element $a_{jk}$ of the product matrix is
$a_{jk}=a_{1;j_{1}k_{1}}a_{2;j_{2}k_{2}}\cdots a_{m;j_{m}k_{m}}.$
Lemma 4.1 [8]. Let $A_{i},i=1,2,\cdots,m$, be a square matrix of order $n_{i}$
and $\xi_{ij},~{}j=1,2,\cdots,n_{i}$ be its eigenvalues. Let
$k_{1},k_{2},\cdots,k_{m}$ be non-negative integers, then the
$n_{1}n_{2}\cdots n_{m}$ eigenvalues of the matrix
$A^{k_{1}}_{1}\otimes\cdots\otimes A^{k_{m}}_{m}$ are $\xi_{j_{1}j_{2}\cdots
j_{m}}=\xi^{k_{1}}_{1j_{1}}\cdots\xi^{k_{m}}_{1j_{m}}$ for $1\leq j_{i}\leq
n_{i}.$
Let $k_{p}=(k_{p1},k_{p1},\cdots,k_{pm}),p=1,2,\cdots,q$ be vectors of non-
negative integers. Then the $n_{1}n_{2}\cdots n_{m}$ eigenvalues of
$\sum\limits_{p=1}^{q}A^{k_{p1}}_{1}\otimes\cdots\otimes A^{k_{pm}}_{m}$ are
$\xi_{j_{1}j_{2}\cdots
j_{m}}=\sum\limits_{p=1}^{q}\xi^{k_{p1}}_{1j_{1}}\cdots\xi^{k_{pm}}_{mj_{m}}.$
For NEPS in graphs see [6, 21]. The following definition extends this concept
to sidigraphs.
Definition 4.2. Let $\mathscr{B}$ be a set of binary $n-$tuples called basis
for the product such that for every $i=1,2,\cdots,n$ there exists
$\beta\in\mathscr{B}$ with $\beta_{i}=1$. The non-complete extended $p-$sum
(or simply called NEPS) of sidigraphs $S_{1},S_{2},\cdots,S_{n}$ with basis
$\bf\it B$ denoted by NEPS$(S_{1},S_{2},\cdots,S_{n};\mathscr{B})$ is a
sidigraph with vertex set $V(S_{1})\times V(S_{2})\times\cdots\times
V(S_{n})$. There is an arc from vertex $(u_{1},u_{2},\cdots,u_{n})$ to
$(v_{1},v_{2},\cdots,v_{n})$ if and only if there exists
$(\beta_{1},\beta_{2},\cdots,\beta_{n})\in\mathscr{B}$ such that
$(u_{i},v_{i})\in A(S_{i})$ whenever $\beta_{i}=1$ and $u_{i}=v_{i}$ whenever
$\beta_{i}=0$. The sign of the arc is given by
$\sigma((u_{1},u_{2},\cdots,u_{n}),(v_{1},v_{2},\cdots,v_{n}))=\prod\limits_{i=1}^{n}\sigma_{i}(u_{i},v_{i})^{\beta_{i}}=\prod\limits_{i:\beta_{i}=1}\sigma_{i}(u_{i},v_{i}).$
Assume that the basis $\mathscr{B}$ has $r\geq 1$ elements, i.e.,
$\beta=\\{\beta_{1},\beta_{2},\cdots,\beta_{r}\\}\subseteq\\{0,1\\}^{n}\backslash\\{(0,0,\cdots,0)\\}$,
we define
NEPS$(S_{1},S_{2},\cdots,S_{n};\mathscr{B})=\bigcup\limits_{\beta\in\mathscr{B}}$
NEPS$(S_{1},S_{2},\cdots,S_{n};\beta).$
Example 4.3. The Kronecker product $S_{1}\otimes S_{2}\otimes\cdots\otimes
S_{n}$ of sidigraphs $S_{1},S_{2},\cdots,S_{n}$ is the NEPS of these
sidigraphs with basis $\beta=\\{(1,1,\cdots,1)\\}$; the cartesian product
$S_{1}\times S_{2}\times\cdots\times S_{n}$ is NEPS with basis
$\\{e_{i}\\},i=1,2,\cdots,n$, where $e_{i}$ is $n-$tuple with $1$ at $i$th
position and $0$ otherwise.
The following result shows that two different basis give disjoint arc sets.
The proof is straightforward.
Lemma 4.4. Let $S=$ NEPS $(S_{1},S_{2},\cdots,S_{n};\beta)$ and $S^{\prime}=$
NEPS $(S_{1},S_{2},\cdots,S_{n};\beta^{\prime})$, $\beta\neq\beta^{\prime}$.
Then $\mathscr{A}(S)\cap\mathscr{A}(S^{\prime})=\phi.$
The following result gives adjacency matrix and spectra of NEPS in terms of
the constituent factor sidigraphs.
Theorem 4.5. If $S=$NEPS$(S_{1},S_{2},\cdots,S_{n};\mathscr{B})$, then the
ajacency matrix is given by
$A(S)=\sum\limits_{\beta\in\mathscr{B}}A^{\beta_{1}}_{1}\otimes\cdots\otimes
A^{\beta_{n}}_{n},$
and eigenvalues are given by $z_{j_{1}j_{2}\cdots
j_{n}}=\sum\limits_{\beta\in\mathscr{B}}z^{\beta_{1}}_{1j_{1}}\cdots
z^{\beta_{n}}_{nj_{n}}$, where $1\leq j_{i}\leq n_{i},i=1,2,\cdots,n.$
Proof. Let $u=(u_{1j_{1}},u_{2j_{2}},\cdots,u_{nj_{n}})$ and
$v=(v_{1j_{1}},v_{2j_{2}},\cdots,v_{nj_{n}})$, where $1\leq
j_{i}\leq|V(S_{i})|$, $i=1,2,\cdots,n$ be any two vertices of $S$. Then
$\displaystyle[A(S)]_{uv}$
$\displaystyle=\sum\limits_{\beta\in\mathscr{B}}(A^{\beta_{1}}_{1})_{u_{1j_{1}}v_{1k_{1}}}(A^{\beta_{2}}_{2})_{u_{2j_{2}}v_{2k_{2}}}\cdots(A^{\beta_{n}}_{n})_{u_{nj_{n}}v_{nk_{n}}}$
$\displaystyle=\sigma_{1}(u_{1j_{1}},v_{1k_{1}})^{\beta_{1}}\sigma_{2}(u_{2j_{2}},v_{2k_{2}})^{\beta_{2}}\cdots\sigma_{n}(u_{nj_{n}},v_{nk_{n}})^{\beta_{n}}$
$\displaystyle=a^{\beta_{1}}_{1;j_{1}k_{1}}a^{\beta_{2}}_{2;j_{2}k_{2}}\cdots
a^{\beta_{n}}_{n;j_{n}k_{n}}=[\sum\limits_{\beta\in\mathscr{B}}A^{\beta_{1}}_{1}\otimes\cdots\otimes
A^{\beta_{n}}_{n}]_{uv}.$
The second part of the result follows by Lemma $4.1$. ∎
We note two special cases of Theorem $4.5$.
$(I)$ The Kronecker product $S_{1}\otimes S_{2}\otimes\cdots\otimes S_{n}$ has
eigenvalues $z_{j_{1}j_{2}\cdots j_{n}}=z_{1j_{1}}z_{2j_{2}}\cdots
z_{nj_{n}}$, $1\leq j_{i}\leq|V(S_{i})|,i=1,2,\cdots,n$.
$(II)$ The cartesian product $S_{1}\times S_{2}\times\cdots\times S_{n}$ has
eigenvalues $z_{j_{1}j_{2}\cdots
j_{n}}=z_{1j_{1}}+z_{2j_{2}}+\cdots+z_{nj_{n}}$, $1\leq
j_{i}\leq|V(S_{i})|,i=1,2,\cdots,n.$
In [8] the authors considered the problem of balance in NEPS of sigraphs. It
is natural to consider the problem of cycle balance for sidigraphs, in view of
Theorem $1.3$, the cartesian product $S_{1}\times S_{2}\times\cdots\times
S_{n}$ of sidigraphs $S_{1},S_{2},\cdots,S_{n}$, is cycle balanced if and only
if $S_{i}$, $i=1,2,\cdots,n$ is cycle balanced. The next result gives
sufficient but not necessary condition for cycle balance of NEPS and the proof
follows on same lines as that in undirected case.
Theorem 4.6. NEPS$(S_{1},S_{2},\cdots,S_{n};\mathscr{B})$ is balanced if
$S_{1},S_{2},\cdots,S_{n}$ are cycle balanced.
Remark 4.7 $(I)$ Theorem $4.6$ does not have a general converse. A counter
example is $S=$ NEPS$(-C_{3},-\overleftrightarrow{K_{2}},\\{(1,1)\\})$, where
$-C_{3}$ denotes all negative directed cycle of order $3$ and
$-\overleftrightarrow{K_{2}}$ is symmetric sidigraph of order $2$ with both
arcs negative. $S$ is all positive, and hence cycle balanced. However $-C_{3}$
is non cycle balanced.
$(II)$ In view of Theorem $1.3$, the converse of Theorem $4.6$ is always true
if basis $\mathscr{B}=\\{e_{i}\\},i=1,2,\cdots,n$.
Theorem 4.8. The following statements are equivalent about cartesian product
$S=S_{1}\times S_{2}\times\cdots\times S_{n}$.
$(I)$ $S$ is cycle balanced.
$(II)$ All of $S_{1},S_{2},\cdots,S_{n}$ are cycle balanced.
$(III)$ $S$ and $S^{u}$ are cospectral.
Proof. Theorem $1.3$ implies equivalence of $(I)$ and $(III)$ and Theorem
$4.6$ implies equivalence of $(I)$ and $(II)$. ∎
## 5\. Upper bounds for the energy of sidigraphs
Let $S$ be a sidigraph of order $n$ with adjacency matrix $A(S)=(a_{ij})$. The
powers of $A(S)$ count the number of walks in signed manner. Let
$w^{+}_{ij}(l)$ and $w^{-}_{ij}(l)$ respectively denote the number of positive
and negative walks of length $l$ from $v_{i}$ to $v_{j}$. The following result
relates integral powers of adjacency matrix with the number of positive and
negative walks.
Theorem 5.1. If $A$ is a adjacency matrix of a sidigraph on $n$ vertices, then
$[A^{l}]_{ij}=w^{+}_{ij}(l)-w^{-}_{ij}(l)$.
Proof. We prove the result by induction on $l$. For $l=1$, result is vacuously
true. For $n=2$, let $n^{+}_{ij}$ denote the number of positive neighbours of
distinct vertices $v_{i}$ and $v_{j}$, $n^{-}_{ij}$ the number of their common
negative neighbours and $n^{\pm}_{ij}$ the number of neighbours that are
positive to one vertex and negative to other. The $(i,i)$ entry of $A^{2}$
equals $w^{+}_{ii}(2)-w^{+}_{ii}(2)$. For $(i,j)$, $i\neq j$,
$n^{+}_{ij}+n^{-}_{ij}=w^{+}_{ij}(2)$ and $n^{\pm}_{ij}=w^{-}_{ij}(2)$, so
that $(i,j)$th entry$=w^{+}_{ij}(2)-w^{-}_{ij}(2)$. Now assume the result to
be true for $l=m$.
We have,
$[A^{m+1}]_{ij}=[A^{m}A]_{ij}=\sum\limits_{k=1}^{n}[A^{m}]_{ik}[A]_{kj}=w^{+}_{ij}(m+1)-w^{-}_{ij}(m+1)$,
by induction hypothesis. Therefore the result follows. ∎
In sidigraph $S$, let $c^{+}_{m}$ denote number of positive closed walks of
length $m$ and $c^{-}_{m}$ the number of negative closed walks of length $m$.
In view of the fact that sum of eigenvalues of a matrix equals to its trace,
we have following Corollary.
Corollary 5.2. If $z_{1},z_{2},\cdots,z_{n}$ are the eigenvalues of a
sidigraph $S$, then $\sum\limits_{j=1}^{n}z^{m}_{j}=c^{+}_{m}-c^{-}_{m}$.
Lemma 5.3. Let $S$ be a sidigraph having $n$ vertices and $a$ arcs and let
$z_{1},z_{2},\cdots,z_{n}$ be its eigenvalues. Then
$(I)$ $\sum\limits_{j=1}^{n}(\Re z_{j})^{2}-\sum\limits_{j=1}^{n}(\Im
z_{j})^{2}=c^{+}_{2}-c^{-}_{2}$, $(II)$ $\sum\limits_{j=1}^{n}(\Re
z_{j})^{2}+\sum\limits_{j=1}^{n}(\Im z_{j})^{2}\leq a=a^{+}+a^{-}$.
Proof. By Corollary $5.2$ , we have
$c^{+}_{2}-c^{-}_{2}=\sum\limits_{j=1}^{n}z_{j}^{2}=\sum\limits_{j=1}^{n}(\Re
z_{j})^{2}-\sum\limits_{j=1}^{n}(\Im z_{j})^{2}+2\iota\sum\limits_{j=1}^{n}\Re
z_{j}\Im z_{j}.$
Equating real and imaginary parts proves $(I)$.
By Schur’s unitary triangularization, there exists a unitary matrix $U$ such
that adjacency matrix $A$ of sidigraph $S$ is unitarily similar to an upper
triangular matrix $T=(t_{jk})$ with $t_{jj}=z_{j}$ for each $j=1,2,\cdots,n$.
Then
$\sum\limits_{j,k=1}^{n}|a_{jk}|^{2}=\sum\limits_{j,k=1}^{n}|t_{jk}|^{2}$. As
$A$ is $(-1,0,1)-$matrix, we have
$a=\sum\limits_{j,k=1}^{n}|\sigma(v_{j},v_{k})|=\sum\limits_{j,k=1}^{n}|a_{jk}|=\sum\limits_{j,k=1}^{n}|a_{jk}|^{2}=\sum\limits_{j,k=1}^{n}|t_{jk}|^{2}\geq\sum\limits_{j=1}^{n}|t_{jj}|^{2}$
$~{}~{}=\sum\limits_{j=1}^{n}|z_{j}|^{2}=\sum\limits_{j=1}^{n}\Re
z_{j}^{2}+\sum\limits_{j=1}^{n}\Im z_{j}^{2}.$
thereby proving $(II)$. ∎
Theorem 5.4. Let $S$ be a sidigraph with $n$ verices and $a=a^{+}+a^{-}$ arcs,
and let $z_{1},z_{2},\cdots,z_{n}$ be its eigenvalues. Then
$E(S)\leq\sqrt{\frac{1}{2}n(a+c^{+}_{2}-c^{-}_{2})}.$ Proof. Subtracting part
$(I)$ of Lemma $5.3$ from $(II)$, we see $\sum\limits_{j=1}^{n}(\Im
z_{j})^{2}\leq\frac{1}{2}(a-(c^{+}_{2}-c^{-}_{2}))$. Applying Cauchy-Schwartz
inequality to vectors $(|\Re z_{1}|,|\Re z_{2}|,\cdots,|\Re z_{n}|)$ and
$(1,1,\cdots,1)$, we have
$E(S)=\sum\limits_{j=1}^{n}|\Re
z_{j}|\leq\sqrt{n}\sqrt{\sum\limits_{j=1}^{n}(\Re
z_{j})^{2}}=\sqrt{n}\sqrt{(c^{+}_{2}-c^{-}_{2})+\sum\limits_{j=1}^{n}(\Im
z_{j})^{2}}$
$\leq\sqrt{n}\sqrt{(c^{+}_{2}-c^{-}_{2})+\frac{1}{2}(a-(c^{+}_{2}-c^{-}_{2})}=\sqrt{\frac{1}{2}n(a+c^{+}_{2}-c^{-}_{2})}.$
∎
Remark 5.5. $(I)$ The upper bound in Theorem $5.4$ is attained by sidigraphs
$S_{1}=(\frac{n}{2}\overleftrightarrow{K_{2}},+)$,
$S_{2}=(\frac{n}{2}\overleftrightarrow{K_{2}},-)$, and skew-smmetric sidigraph
of order $n$. Note that eigenvalues of $S_{1}$ and $S_{2}$ are $-1,+1$ each
repeated $\frac{n}{2}$ times and eigenvalues of skew symmetric sidigraph of
order $n$ are $0^{n-2},~{}\pm\iota\sqrt{n-1}$.
$(II)$. Above result extends McClleland’s inequality for sigraphs [9] which
states that $E(\Sigma)\leq\sqrt{2pq}$, holds for every sigraph with $p$
vertices and $q$ edges. Let $\overleftrightarrow{\Sigma}$ be the symmetric
sidigraph of sigraph $\Sigma$, then in $\overleftrightarrow{\Sigma}$,
$a=2q=c^{+}_{2}=c^{+}_{2}-c^{-}_{2}$,. By Theorem $5.4$,
$E(\Sigma)=E(\overleftrightarrow{\Sigma})\leq\sqrt{\frac{1}{2}p(2q+2q)}=\sqrt{2pq}.$
The following result gives sharp upper bound of energy of sidigraphs in terms
of number of arcs. The proof is same as in unsigned case.
Theorem 5.6. Let $S$ be a sidigraph with $a$ arcs. Then $E(S)\leq a$ with
equality if and only if $S=(\frac{a}{2}\overleftrightarrow{K_{2}},+)$ or
$S=(\frac{a}{2}\overleftrightarrow{K_{2}},-)$ plus some isolated vertices.
Remark 5.7. Theorem $5.6$ extends the result for sigraphs [9], which states
that $E(\Sigma)\leq 2q$ for every sigraph with $q$ edges with equality if and
only if $\Sigma=(\frac{q}{2}K_{2},+)$ or $\Sigma=(\frac{q}{2}K_{2},-)$ plus
some isolated vertices.
## 6\. Equienergetic sidigraphs
Two sidigraphs are said to be isomorphic if their underlying digraphs are
isomorphic such that signs are preserved. Any two isomorphic sidigraphs are
obviously cospectral. There exist nonisomorphic sidigraphs which are
cospectral, e.g., consider the sidigraphs $S_{1}$ and $S_{2}$ shown in Figure
$2$. Clearly, $S_{1}$ and $S_{2}$ are nonisomorphic.
$Spec~{}S_{1}=\\{0^{(5)}\\}=Spec~{}S_{2}$.
$S_{1}$Figure 2Figure 2$S_{2}$
Two nonisomorphic sidigraphs $S_{1}$ and $S_{2}$ of same order are said to be
equienergetic if $E(S_{1})=E(S_{2})$. In [20] Rada proved the existence pairs
of non-symmetric equienergetic digraphs on $n\geq 3$ vertices. Cospectral
sidigraphs are obviously equienergetic, therefore problem of equienergetic
sidigraphs reduces to problem of construction of noncospectral pairs of
equinergetic sidigraphs such that for every pair not both sidigraphs are cycle
balanced. The problem of construction of equienergetic sigraphs is an open
problem [7] and for equienergetic graphs see [4, 19].
Theorem 6.1. Let $S$ be a sidigraph of order $n$ having eigenvalues
$z_{1},z_{2},\cdots,z_{n}$ such that $|\Re z_{j}|\leq 1$ for every
$j=1,2,\cdots,n$. Then $E(S\times\overleftrightarrow{K_{2}})=2n$.
Proof. Let $z_{1},z_{2},\cdots,z_{t}$ be eigenvalues with nonnegative real
part and $z_{t+1},\cdots,z_{n}$ be those with negative real part. Eigenvalues
of cartesian product $S\times\overleftrightarrow{k_{2}}$ are $z_{1}\pm
1,z_{2}\pm 1,\cdots,z_{t}\pm 1,z_{t+1}\pm 1,\cdots,z_{n}\pm 1$. Therefore
$E(S\times\overleftrightarrow{K_{2}})=\sum\limits_{j=1}^{t}(|\Re z_{j}+1|+|\Re
z_{j}-1|)+\sum\limits_{j=t+1}^{n}(|\Re z_{j}+1|+|\Re z_{j}-1|).$
As $|\Re z_{j}|\leq 1$ for all $i=1,2,\cdots,n$, it follows that
$E(S\times\overleftrightarrow{K_{2}})=\sum\limits_{j=1}^{t}(\Re z_{j}+1+1-\Re
z_{j})+\sum\limits_{j=t+1}^{n}(1-\Re z_{j}+\Re z_{j}+1)=2t+2(n-t)=2n.$
∎
Corollary 6.2. For $n\geq 2$, $E({\bf
C}_{n}\times\overleftrightarrow{K_{2}})=E(C_{n}\times\overleftrightarrow{K_{2}})=2n$.
Moreover, ${\bf C}_{n}\times\overleftrightarrow{K_{2}}$ and
$C_{n}\times\overleftrightarrow{K_{2}}$ are noncospectral sidigraphs with $2n$
vertices.
Proof. We know that eigenvalues of ${\bf C}_{n}$ are
$e^{\frac{\iota(2j+1)\pi}{n}},~{}j=0,1,\cdots,n-1$ and those of $C_{n}$ are
$e^{\frac{2\iota j\pi}{n}},~{}j=0,1,\cdots,n-1$. Clearly, eigenvalues of ${\bf
C}_{n}$ and ${C_{n}}$ meet requirement of Theorem $6.1$, so $E({\bf
C}_{n}\times\overleftrightarrow{K_{2}})=E(C_{n}\times\overleftrightarrow{K_{2}})=2n$.
Moreover, $2\notin spec({\bf C}_{n}\times\overleftrightarrow{K_{2}})$ but
$2\in spec(C_{n}\times\overleftrightarrow{K_{2}})$ implying that ${\bf
C}_{n}\times\overleftrightarrow{K_{2}}$ and
$C_{n}\times\overleftrightarrow{K_{2}}$ are non cospectral. Number of vertices
in both sidigraphs is $2n$ follows by definition of cartesian product. In view
of Theorem $1.3$, ${\bf C}_{n}\times\overleftrightarrow{K_{2}}$ is non cycle
balanced whereas $C_{n}\times\overleftrightarrow{K_{2}}$ is cycle balanced. ∎
Example 6.3. For each odd $n$, ${\bf C}_{n}$ and $C_{n}$ is a non co-spectral
pair of equienergetic sidigraphs, because $Spec({\bf C}_{n})=-Spec(C_{n})$ and
$1\notin spec({\bf C}_{n})$ but $1\in spec(C_{n})$.
From Corollary $6.2$ and Example $6.3$, we see that for each positive integer
$n\geq 3$ there exits a pair of non co-spectral sidigraphs with one sidigraph
cycle balanced and another non cycle balanced. We now construct pairs of non
cospectral sidigraphs of order $2n,~{}n\geq 5$ with both constituents non
cycle balanced. Let $P^{l}_{n}~{}(n\geq l+1)$ be a sidigraph obtained by
identifying one pendant vertex of the path $P_{n-l+1}$ with any vertex of
${\bf C}_{l}$. Sign of non cyclic arcs is immaterial.
Corollary 6.4. For each $n\geq 5$, $P^{3}_{n}\times\overleftrightarrow{K_{2}}$
and $P^{4}_{n}\times\overleftrightarrow{K_{2}}$ is a pair of noncospectral
equienergetic sidigraphs of order and energy equal to $2n$.
Proof. Using the fact that $\phi_{P^{l}_{n}}(x)=x^{n-l}\phi_{{\bf C}_{l}}(x)$
and Theorem $6.1$, it follows that
$E(P^{3}_{n}\times\overleftrightarrow{K_{2}})=E(P^{4}_{n}\times\overleftrightarrow{K_{2}})=2n$.
Now $1$ is an eigenvalue of $P^{3}_{n}\times\overleftrightarrow{K_{2}}$ with
multiplicity $n-3$ but $1$ is eigenvalue of
$P^{4}_{n}\times\overleftrightarrow{K_{2}}$ with multiplicity $n-4$, therefore
these two sidigraphs are noncospectral. Order of both sidigraphs equals to
$2n$ follows by definition of cartesian product. In view of Remark
$4.7~{}(II)$, it follows that both $P^{3}_{n}\times\overleftrightarrow{K_{2}}$
and $P^{4}_{n}\times\overleftrightarrow{K_{2}}$ are non cycle balanced. ∎
Corollary 6.5. If $S$ is a sidigraph on $n$ vertices having eigenvalues
$z_{1},z_{2},\cdots,z_{n}$ such that $|\Im z_{j}|\leq\frac{1}{\sqrt{n-1}}$ for
all $j=1,2,\cdots,n$, then $E((S\otimes{\bf
S}_{m})\times\overleftrightarrow{K_{2}})=2nm=$ order of sidigraph
$(S\otimes{\bf S}_{m})\times\overleftrightarrow{K_{2}}$, $2\leq m\leq n$.
Proof. Eigenvalues of $S$ are $z_{1},z_{2},\cdots,z_{n}$ with $|\Im
z_{j}|\leq\frac{1}{\sqrt{n-1}}$ and eigenvalues of ${\bf S}_{m}$ are
$0^{m-2},\pm\iota\sqrt{m-1}$. Then eigenvalues of Kronecker product
$S\otimes{\bf S}_{m}$ meet requirement of Theorem $6.1$, therefore
$E((S\otimes{\bf S}_{m})\times\overleftrightarrow{K_{2}})=2nm$. Order of
sidigraph equals to $2nm$ follows from definition of Kronecker product and
cartesian product. ∎
Theorem 6.6. Let $A=(a_{ij})$ be a square matrix of order $n$ having integral
entries and zero trace and let $z_{1},z_{2},\cdots,z_{n}$ be its eigenvalues.
Put $\alpha=\sum\limits_{j=1}^{n}|\Re z_{j}|$, then $\alpha$ cannot be of the
form $(I)~{}(2^{t}s)^{\frac{1}{h}}$ with $h\geq 1,~{}0\leq t<h$ and $s$ odd
$(II)~{}(\frac{m}{n})^{\frac{1}{r}}$ where $\frac{m}{n}$ is non-integral
rational.
Proof. We note that $\alpha=2\sum\limits_{\Re z_{j}\geq 0}z_{j}$. Put
$z=\sum\limits_{\Re z_{j}\geq 0}z_{j}$, then $z$ being sum of algebraic
integers is an algebraic integer.
$(I)$ Assume $\alpha=(2^{t}s)^{\frac{1}{h}}$, so that
$2z=(2^{t}s)^{\frac{1}{h}}$. Simplifying gives $z^{h}=\frac{s}{2^{l}}$, where
$l=h-t\geq 1$. As $s$ is odd, therefore we see $z^{h}$ is non-integral
rational algebraic integer, a contradiction.
$(II)$ Proof is same as in part $(I)$. ∎
Bapat and Pati [3] characterized positive rationals which can be the energy of
graph. Later Pirzada and Gutman [16] proved energy of graph cannot be square
root of an odd integer. These results were generalized to digraphs by Pirzada,
Mushtaq, Gutman and Rada [17]. Following result generalizes these results to
sidigraphs.
Theorem 6.7. Energy of a sidigraph cannot be of the form
$(I)~{}(2^{t}s)^{\frac{1}{h}}$ with $h\geq 1,~{}0\leq t<h$ and $s$ odd
$(II)~{}(\frac{m}{n})^{\frac{1}{r}}$ where $\frac{m}{n}$ is non-integral
rational number and $r\geq 1$.
Proof. Let $S$ be a sidigraph with adjacency matrix $A(S)$. Apply Theorem
$6.6$ to $A(S)$ and note that $E(S)=\sum\limits_{j=1}^{n}|\Re z_{j}|$, the
result follows. ∎
## 7\. Open problems
We conclude this paper with the following open problems.
$(1)$ Characterize sidigraphs with energy equal to the number of vertices.
$(2)$ Find an infinite family of noncospectral equienergetic sidigraphs on
$n\geq 4$ vertices with both constituents non cycle balanced.
$(3)$ Determine bases other than $\\{e_{i}\\}$ for which converse of Theorem
$4.7$ holds.
Acknowledgements The second author thanks University Grants Commission, New
Delhi, India for providing junior research fellowship.
## References
* [1] B. D. Acharya, Spectral criterion for the cycle balance in networks. J. Graph Theory 4(1980)1-11.
* [2] B. D. Acharya, M. K. Gill, G. A. Patwardhan, Quasicospectral graphs and digraphs, National symposium on mathematical modelling M. R. I. Allahabad: July 19-20, 1982.
* [3] R. B. Bapat, S. Pati, Energy of a graph is never an odd integer, Bull. Kerala Math. Assoc., 1(2004), 129-132.
* [4] A. S. Bonifaco, C. T. M. Vinagre, N. M. M. Abreu, Constructing pairs of equienergetic and non cospectral graphs, Applied Mathematics Letters 21(2008) 338–341.
* [5] R. Brualdi, Spectra of digraphs, Linear Algebra and Appl., 432(2009), 2181-2213.
* [6] D. M. Cvetkovic, M. Doob and H. Sachs, Spectra Of Graphs. Academic press, New York 1980.
* [7] K. A. Germina and Shahul Hameed K, On signed paths, signed cycles and their energies. Appl. Math. Sci. 4(2010), no. 70, 3455-3466.
* [8] K. A. Germina and Shahul Hameed K, and Thomas Zaslavsky, On products and line graphs of signed graphs, their eigenvalues and energy, Linear Algebra and its Applications, 435, 10 (2010) 2432-2450.
* [9] K. A. Germina and Shahul Hameed K, On bounds of eigenvalues and energy of signed graphs: submitted.
* [10] M. K. Gill and B. D. Acharya, A recurrence formula for computing the characteristic polynomial of a sigraph, Journal of Combinatorics, information and system sciences, 5(1980) 68-72.
* [11] I. Gutman, The Energy of a graph, Ber. Math. Statist. Sekt. Forschungszenturm Graz., 103(1978), 1-22.
* [12] I. Gutman, The Energy of a graph: Old and new results, in: A. Betten, A. Kohnert, R. Laue, A. Wasserman (Eds), Algebriac Combinatorics and Applications, Springer-Verlag, Berlin 2001, pp. 196-211.
* [13] R. Horn and C. Johnson, Matrix Analysis; Cambridge University press, 1985.
* [14] W. Lopez and J. Rada, Equienergetic digraphs, International journal of pure and applied mathematics, 36(2007) 361-372.
* [15] I. Pena, J. Rada, Energy of digraphs, Linear and Multilinear Algebra, 56(5) (2008), 565-579.
* [16] S. Pirzada and I. Gutman, Energy of a graph is never the square root of an odd integer, Appl. Anal. Dis. Math., 2(2008), 118-121.
* [17] S. Pirzada, Mushtaq A. Bhat, I. Gutman, J. Rada, On the energy of digraphs, Bulletin of International Mathematical Virtual Institute, 3(2013), 69-76.
* [18] S. Ponnusamy, Foundations of complex analysis (Narosa).
* [19] H. S. Ramane, H. B. Walikar, Construction of equienergetic graphs, MATCH Commun. Math. Comput. Chem., 57(2007) 23-210.
* [20] J. Rada, The McClelland inequality for the energy of digraphs, Linear Algebra and its Applications, 430(2009), 800-804.
* [21] D. Stevanoic, Energy and NEPS of graphs, Linear and Multilinear Algebra, 53(2005) 67-74.
* [22] T. Zaslavsky, Matrices in the theory of signed simple graphs. In: B. D. Acharya, G. O. H. Katona, and J. Nesetril, eds., Advances in Discrete Mathematics and Applications: Mysore, 2008(ICDM-2008, Mysore, India), pp. 207-229. Ramanajun Math. Soc. Lect. Notes Ser., No. 13. Ramanajun Mathematical Soc., Mysore, India, 2010.
|
arxiv-papers
| 2013-09-24T17:50:58 |
2024-09-04T02:49:51.455525
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "S. Pirzada and Mushtaq A. Bhat",
"submitter": "Shariefuddin Pirzada",
"url": "https://arxiv.org/abs/1309.6266"
}
|
1309.6286
|
# Controllable emission of a dipolar source coupled with a magneto-dielectric
resonant subwavelength scatterer
Brice Rolly1, Jean-Michel Geffrin1, Redha Abdeddaim1, Brian Stout1, Nicolas
Bonod1∗
1 CNRS, Aix-Marseille Université, Centrale Marseille, Institut Fresnel, UMR
7249, Campus de St Jérôme, 13397 Marseille, France
∗ Corresponding author: [email protected]
We demonstrate experimentally and theoretically that a local excitation of a
single scatterer of relative dielectric permittivity $\varepsilon=6$ permits
to excite broad dipolar and quadrupolar electric and magnetic resonances that
shape the emission pattern in an unprecedented way. By suitably positioning
the feed with respect to the sphere at a $\lambda/3$ distance, this compact
antenna is able to spectrally sort the electromagnetic emission either in the
forward or in the backward direction, together with a high gain in
directivity. Materials with $\varepsilon=6$ can be found in the whole spectrum
of frequencies promising Mie antennas to become an enabling technology in
numbers of applications, ranging from quantum single photon sources to
telecommunications.
## Introduction
Mie resonances in spherical particles of moderate refractive index were
recently shown to exhibit well pronounced electric and magnetic modes of both
dipole and quadrupole orders [1, 2, 3, 4, 5, 6, 7, 8]. As first predicted by
Kerker in the case of particles with magnetic permeability [9], interference
effects between electric and magnetic modes can strongly favour either forward
or backward scattering, respectively depending on the frequency of the
incident plane wave [10, 11, 12, 13, 14, 15] with applications in anti-
reflecting structured surfaces for photovoltaic cells [16]. Dielectric
particles reveal to be highly interesting scatterers to design highly
directive, compact and lossless antennas [17, 18, 19, 20]. Recently, it was
proposed to extend the so-called Kerker’s conditions to the case of near field
excitations, where the interplay between electric and magnetic modes can boost
the directivity of antennas [21]. It was also proposed to realize a notch
inside a dielectric scatterer of refractive index $n$ around 4 in order to
increase the number of multipoles excited by a near field source [26]. Such
magneto-dielectric antennas offer more parameters for optimizing the emission
directivity than classical metallic elements which behave principally as
electric multipolar scatterers [27, 28, 29, 30, 31, 32] and they could be very
interesting to direct the fluorescence emission of quantum emitters in the
optical range of frequencies.
Here we propose to design a versatile antenna with the use of a single
spherical scatterer exhibiting a refractive index $n=2.45$. Mie resonators
with refractive index typically around 3.5 have been mainly considered because
they exhibit narrow and well pronounced electric and magnetic modes. A
decrease of the refractive index leads to a broadening of the Mie resonances
and a refractive index $n=2.45$ allows for a remarkable mixing between
electric and magnetic dipolar and quadrupolar modes over a finite range of
frequencies. We show that the interplay between the first two electric and
magnetic modes permits to control the emission directivity by tuning the
emission frequency. Also, lessening the refractive index opens the way towards
the use of oxides such as rutile TiO2 to design highly directive optical
antennas.
## Results
We measure the emission pattern in the E-plane (i.e. $xOz$ plane, see Fig.1d)
of a subwavelength scale ‘Mie’ antenna consisting of a single dielectric
particle of permittivity $\varepsilon=6$ ($n$=2.45) fed in the GHz regime by a
two-arm electric dipole emitter. By controlling the emitter-to-particle
distance at a deep subwavelength scale, we report on the possibility to choose
the emission direction by tuning the frequency from 8.7 GHz to 9.74 GHz.
Experiments are carried out in an anechoic chamber (CCRM-Marseille, see
Fig.1a) dedicated to amplitude and phase measurements of the electric field
with a receiver antenna rotating along a circular arm 4 m in diameter,
centered on the feed [13]. Experimental measurements are accompanied by a
theoretical derivation of the emission pattern based on the coupling between
the feed and the magnetic and electric modes of the sphere.
Figure 1: (a) Photograph of the experiment in the anechoic chamber: the
dipole-sphere antenna is placed on a polystyrene mast at the center of an
anechoic chamber. The receiver antenna (ARA DRG 118) can rotate around the
$Oy$ axis (see d) except in an exclusion zone in the angle range
$[-130^{\circ},130^{\circ}]$ due to the presence of the vertical arch. (b)
Close-up of the dipole-sphere antenna and of the polystyrene holders. (c) Norm
of the electric field emitted by the isolated dipole antenna as a function of
the frequency and angle $\phi_{\rm exp}$ in the E-plane ($xOz$ plane) in dB.
At each frequency, the field is normalized by its maximum in the E-plane. (d)
Sketch of the experiment with the axes and $\phi_{\rm exp}$ angle in the
$(xOz)$ plane, used in Fig.1c, Fig.2 and Fig.3; the blue arrow indicates the
direction of the receiving antenna. (e) Norm of the Mie coefficients
$|c_{j}^{\mathrm{e(m)}}|$ with respect to frequency in GHz. Black line:
magnetic dipole $|c_{1}^{\mathrm{m}}|$; red line: electric dipole
$|c_{1}^{\mathrm{e}}|$; green line: magnetic quadrupole
$|c_{2}^{\mathrm{m}}|$; blue line: electric quadrupole $|c_{2}^{\mathrm{e}}|$.
The two-wire dipole is directly attached to a connector (SubMiniature version
A). Each arm is 9 mm long which results in a total length smaller than the
wavelength even at the highest operating frequency. The characterization of
the dipole in its E-plane (Fig. 1c) shows a classical dipole-like radiation
pattern, with comparable front and backward radiated amplitude (respectively
$0^{\circ}$ and $180^{\circ}$), and negligible emission in the $90^{\circ}$
direction. The emitter is coupled with a single dielectric sphere, 19 mm in
diameter, composed of an Eccostock HIK (Emmerson $\&$ Cumming) material of
permittivity 6, presenting low losses ($\tan\delta\leq 0.02$). The distance
between the dipole and the sphere can be controlled at a submillimeter scale
via an expanded polystyrene holder equivalent to air at the operating
frequencies (see Fig.1a-b). Measurements have been performed with emitter-to-
sphere distances of 5, 10 and 20 mm and frequencies ranging from 8 to 10 GHz
with a step of 20 MHz, which corresponds to a total of 300 emission pattern
measurements (the emitter-to-sphere distance is defined between the emitter
and the surface of the sphere, i.e. it is equal to emitter-to-sphere center
distance minus the sphere radius). The amplitude of the electric field is
displayed in Fig. 2 with respect to the frequency and angle $\phi_{\rm exp}$
ranging from 0 to 180∘, i.e. in a half-plane containing one of the emitter
arms. At every frequency, the amplitude is normalized by the maximum measured
in the [0;180∘] range, and is displayed in dB. One observes that the radiation
pattern is highly sensitive to modifications of the emitter-to-particle
distance on a scale much smaller than the emission wavelength. Importantly, we
observe that this Mie antenna can emit either in the forward (around 8.75 GHz)
or backward (around 9.5 GHz) direction efficiently at a 10 mm distance.
In order to demonstrate the contribution of each mode in tuning the emission
directivity, we derived a formula of the irradiance in the E-plane (derivation
is detailed in the Supplementary information). Such a formula invokes the far
field expression of the interference between the fields produced by (i) the
dipolar source and (ii) the induced electric-magnetic dipoles and quadrupoles
in the dielectric particle. The fields are derived in the spherical vector
basis,
$[\hat{\mathbf{e}}_{\mathrm{r}},\hat{\mathbf{e}}_{\mathrm{\theta}},\hat{\mathbf{e}}_{\mathrm{\phi}}]$
with the dipole emitter defined to lie along the $+z$ axis
($\mathbf{p}_{0}\equiv p_{0}\hat{\mathbf{z}}$) and placed at the origin of the
coordinate system. The center of the spherical scatterer is placed at a
distance $d$ from the emitter on the $x$-axis,
$\mathbf{r}_{1}=+d\hat{\mathbf{x}}$. The total, normalized (respectively to
the electric dipole emitter maximum far-field irradiance) far-field irradiance
can be cast in the standard spherical coordinate system:
$\displaystyle I(\theta,\phi)$ $\displaystyle=$
$\displaystyle|\tilde{\mathbf{E}}(\mathbf{r})|^{2}$ (1) $\displaystyle=$
$\displaystyle\bigg{|}\sin\theta\hat{\mathbf{e}_{\theta}}$
$\displaystyle+e^{i\varphi}\gamma^{e}_{1}\tilde{\alpha}^{e}_{1}\sin\theta\hat{\mathbf{e}_{\theta}}$
$\displaystyle+e^{i\varphi}\gamma^{m}_{1}\tilde{\alpha}^{m}_{1}(\cos\phi\hat{\mathbf{e}_{\theta}}-\sin\phi\cos\theta\hat{\mathbf{e}_{\phi}})$
$\displaystyle+e^{i\varphi}\gamma^{e}_{2}\tilde{\alpha}^{e}_{2}(\cos\phi\cos
2\theta\hat{\mathbf{e}_{\theta}}-\sin\phi\cos\theta\hat{\mathbf{e}_{\phi}})$
$\displaystyle+e^{i\varphi}\gamma^{m}_{2}\tilde{\alpha}^{m}_{2}(\cos
2\phi\sin\theta\hat{\mathbf{e}_{\theta}}-\frac{\sin 2\phi\sin
2\theta}{2}\hat{\mathbf{e}_{\phi}})\bigg{|}^{2}$
In the latter expression, each subsequent line stands for the field produced
by the source and the induced electric dipole, magnetic dipole, electric
quadrupole, and magnetic quadrupole respectively; $e^{i\varphi}$ is a far-
field phase shift contribution; $\tilde{\alpha}_{n}^{e(m)}$ is a
_dimensionless_ electric (magnetic) polarizability of order $n$, defined in
terms of Mie theory and a T-matrix formalism (see Supporting Information). The
first four Mie coefficients are plotted in Fig.1e with respect to the
frequency and one remarks that the relative low refractive index of the
resonant scatterer leads to broad resonances and that the norms of the four
Mie coefficients are nearly equal near 9.75 GHz. The coupling coefficients
between the emitter and the first 2 electric and magnetic modes of the sphere
can be cast [33]:
$\displaystyle\gamma^{e}_{1}$ $\displaystyle\equiv$ $\displaystyle
e^{ikd}(k^{2}d^{2}+ikd-1)(a/d)^{3}$ $\displaystyle\gamma^{m}_{1}$
$\displaystyle\equiv$ $\displaystyle e^{ikd}(k^{2}d^{2}+ikd)(a/d)^{3}$
$\displaystyle\gamma^{e}_{2}$ $\displaystyle\equiv$
$\displaystyle-\frac{5}{3}e^{ikd}(k^{3}d^{3}+3ik^{2}d^{2}-6kd-6i)(a/d)^{4}$
$\displaystyle\gamma^{m}_{2}$ $\displaystyle\equiv$
$\displaystyle\frac{5}{3}e^{ikd}(k^{3}d^{3}+3ik^{2}d^{2}-3kd)(a/d)^{4}$
In the $(xOz)$ ($\phi=0$ or $\pi$) plane, the irradiance can be cast:
$\displaystyle I^{\pm xz}(\theta)$ $\displaystyle=$
$\displaystyle\bigg{|}\sin\theta(1+e^{-ikd\pm\sin\theta}\gamma^{e}_{1}\tilde{\alpha}^{e}_{1})$
(2)
$\displaystyle-e^{-ikd\pm\sin\theta}\times\big{(}\pm\gamma^{m}_{1}\tilde{\alpha}^{m}_{1}\pm\gamma^{e}_{2}\tilde{\alpha}^{e}_{2}\cos
2\theta$
$\displaystyle+\gamma^{m}_{2}\tilde{\alpha}^{m}_{2}\sin\theta\bigg{|}^{2}$
Figure 2: Norm of the electric field in the E-plane, in decibels, measured
(top row) and simulated (bottom row) as functions of the antenna receiver
angle $\phi_{\rm exp}$ in degrees (in abscissa) and emitting frequency in GHz
(in ordinate). At each frequency, the field is normalized by its maximum in
the E-plane. The emitter-to-sphere gaps are (a) 5 mm, (b) 10 mm and (c) 20 mm.
The normalized amplitude of the electric field in the E-plane is displayed in
Fig.2 as a function of the angle in the $(xOz)$ plane, and a good agreement
between theory and experiments is observed for the three emitter-to-particle
distances despite a small frequency shift that is likely due to an imperfect
knowledge of the permittivity of the sphere (the real part is given at
$\pm{5}\%$ and the imaginary part is set to 0 in the model) and to the
physical length of the source, the theoretical expression assuming a point
dipole. This model is general, and could also be applied to predict the
scattering pattern of metallic magneto-electric antennas provided that the
diagonal polarizability tensor elements are accurately determined and
predominate over the extra-diagonal elements.
## Discussion
Both theoretical model and experiments predict that the privileged direction
of emission of the antenna can be controlled for a 10 mm separation distance
by tuning the emission frequency. This feature is highlighted in Figs.3a and b
where the intensities of the scattered electric field measured in the E-plane
at frequencies 8.7 GHz and 9.74 GHz are plotted in polar coordinates.
Measurements and theory are quantified by defining directivity gain, in
isotropic decibels, $D_{\mathrm{dBi}}=10\log\big{(}4\pi
I/P_{\mathrm{rad}}\big{)}$ where $I$ and $P_{\mathrm{rad}}$ are respectively
the radiant intensity in the direction of interest and the total radiative
power of the antenna. Theoretical values of 7.01 dBi at 8.66 GHz and 5.17 dBi
at 9.58 GHz in the forward and backward directions respectively can be
achieved. The sphere also effectively increases antenna gain: the presence of
the sphere increases the measured electric field intensity: by a factor of 3.4
and 2.4 along the corresponding privileged directions at 8.7 GHz and 9.74 GHz
respectively.
Figure 3: (a,b) Norm of the electric field scattered for an emitter-to-
particle gap of 10 mm in the E-plane, (red line) simulations with 20
multipoles and (blue) measurements: (a) Back-scattering observed at 9.74 GHz,
(b) Forward scattering observed at 8.7 GHz. (c) Emission pattern derived with
Eq. 2 when considering the induced dipoles (red line) or quadrupoles (blue
line) only at 8.7 GHz. (d) Emission pattern obtained with Eq. 2 when
considering electric (red line) or magnetic modes (blue line) of dipolar and
quadrupolar orders, sum of electric and magnetic response of both dipole and
quadrupole order (green line), electric and magnetic modes with 20 multipoles
(dashed black line) at 8.7 GHz.
The emission patterns plotted in Fig.3c at $f$=8.7 GHz when considering
dipoles or quadrupoles only (full red and blue lines) show that neither dipole
nor quadrupole excitations taken alone suffice to explain the directivity: a
purely response of the sphere would predict light preferentially emitted in
the backward direction, whereas a solely quadrupole response would lead to
forward emission together with a low gain in directivity. A similar filtering
with respect to either purely electric or magnetic modes, illustrated in
Fig.3d, shows that neither purely electric nor purely magnetic modes can fit
the emission pattern measured in the E-plane, while the emission pattern is
accurately reproduced when both dipoles and quadrupoles are considered.
Comparison with a multipolar calculation performed with a Generalized Mie
theory calculation taking into account 20 multipole orders reveals the high
accuracy of the quadrupolar model. These last two plots demonstrate that the
emission pattern results from an efficient coupling between the electric and
magnetic modes, of both dipolar and quadrupolar orders which is made possible
by the broadness of the Mie resonances displayed in Fig.1e. The high
sensitivity of the emission pattern with respect to the emission frequency is
explained by a strong modulation of three of the relevant modes between 8.5
GHz and 10 GHz (see Fig.1e).
In conclusion, a single dielectric particle of refractive index $n=2.45$ is
used to efficiently tune the scattering directivity of an electric dipole
emitter. The tuning of the emission frequency modulates the amplitude and
phase of the first four modes of the sphere which drastically modifies the
emission pattern. Dielectric particles of refractive index around 2.5 can be
found over a broad spectrum making dielectric Mie antennas interesting from
microwaves to nano-optics where they could acts as lossless subwavelength
spectral sorters [34, 35].
## Methods
### Measurements
The measurements are carried out in an anechoic chamber (Centre Commun de
Ressources en Microondes in Marseille, France) that allows electric field
amplitude and phase measurements thanks to an antenna receiver placed on a
circular arm of radius 2 m. For this study, the electric field is quantified
by measuring the electric field component lying inside the E-plane (the
electric field component normal to the E-plane being negligible). The $S_{21}$
parameter is measured to provide the norm of the electric field $20\log|E|$.
The field is normalized at each frequency by the maximum value $E_{max}$
measured when varying the angle $\phi_{\rm exp}$ and the value of
$20\log|E/E_{max}|$ is reported in Fig. 2. In Fig. 3, the field is also
normalized by its maximum when varying the angle $\phi_{\rm exp}$ at
frequencies 8.7 GHz and 9.74 GHz.
The antenna receiver cannot make a full rotation around the axis due to the
presence of the vertical arch. The exclusion zone is almost $100^{\circ}$. The
complete $360^{\circ}$ emission pattern requires a second measurement for
which the mast holder is rotated by $120^{\circ}$. In Figs.3a and b, it can be
observed that the two measurements match and permit the reconstruction of the
emission patterns with a high accuracy.
### Theoretical calculations
The far field expressions of the electric fields radiated by the feeding
electric dipole, and the electric and magnetic dipoles and quadrupoles excited
in the particle are derived and summed up to obtain the final expression of
the total electric field. The sources are strongly coupled in near field and
the calculation of the phase differences between the coherent fields must take
into account the entire field, i.e. near, intermediate and far field
expressions. The expression of the emission pattern requires the calculation
of the dipolar and quadrupolar polarizabilities of the scatterer, that can be
easily obtained within the context of generalized Mie theory [36]. The
analytical derivation is detailed in the Supporting Information. Numerical
simulations are performed using an in-house Generalized Mie Theory code. The
embedding medium is air and its refractive index is taken equal to 1, the
dielectric permittivity of the spherical scatterer is estimated by Emmerson
Cumming to be equal to $6\pm{5\%}$.
## References
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* [30] Bonod, N., Devilez, A., Rolly, B., Bidault, S. & Stout, B. Ultracompact and unidirectional metallic antennas. _Physical Review B: Condensed Matter and Materials Physics_ 82, 115429 (2010).
* [31] Rolly, B., Stout, B., Bidault, S. & Bonod, N. Crucial role of the emitter–particle distance on the directivity of optical antennas. _Optics Letters_ 36, 3368–3370 (2011).
* [32] Massa, E., Maier, S. A. & Giannini, V. An analytical approach to light scattering from small cubic and rectangular cuboidal nanoantennas. _New Journal of Physics_ 15, 063013 (2013).
* [33] Rolly, B., Stout, B. & Bonod, N. Metallic dimers: When bonding transverse modes shine light. _Physical Review B: Condensed Matter and Materials Physics_ 84, 125420 (2011).
* [34] Aouani, H. _et al._ Bright Unidirectional Fluorescence Emission of Molecules in a Nanoaperture with Plasmonic Corrugations. _Nano Letters_ 11, 637–644 (2011).
* [35] Shegai, T. _et al._ A bimetallic nanoantenna for directional colour routing. _Nature Communications_ 2, 481 (2011).
* [36] Stout, B., Devilez, A., Rolly, B. & Bonod, N. Multipole methods for nanoantennas design: applications to Yagi-Uda configurations. _Journal of the Optical Society of America B: Optical Physics_ 28, 1213–1223 (2011).
## Acknowledgements
The authors thank Sébastien Bidault for discussions and Jean-Pierre Spinelli
for his involvement in the anechoic chambers. This work is supported by Agence
Nationale de la Recherche via project ANR 11 BS10 002 02 TWINS.
## Author contributions
N.B, R.A. and J.M.G. conceived the experiment; R. A. and J.M.G. carried out
the experiment; B. R. derived the analytical expression, performed the
numerical calculations with B. S. and N. B.; All the authors contributed to
the discussion and analysis of the results; N. B. wrote the manuscript with
contribution from all the authors.
## Additional information
The authors declare no competing financial interests. Supplementary
information accompanies this paper and include the detailed derivation of
expression 2. Reprints and permissions information is available at
www.nature.com/reprints. Correspondence and requests for materials should be
addressed to N. B.
|
arxiv-papers
| 2013-09-22T10:56:20 |
2024-09-04T02:49:51.466704
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Brice Rolly, Redha Abdeddaim, Jean-Michel Geffrin, Brian Stout,\n Nicolas Bonod",
"submitter": "Nicolas Bonod",
"url": "https://arxiv.org/abs/1309.6286"
}
|
1309.6386
|
# Flashing coherently rotating carbon sticks in 24Mg+24Mg collision
M. H. Zhao S. Kun [email protected] O. Merlo M. R. Huang Y. Li J. S. Wang
Institute of Modern Physics, Chinese Academy of Sciences, Lanzhou, 730000,
People’s Republic of China University of Chinese Academy of Sciences, Beijing
100049, People’s Republic of China Canberra, Australia Zurich University of
Applied Sciences, Institute of Applied Simulation Grüental, P.O. Box CH-8820
Waedenswil, Schweiz
###### Abstract
We analyze quasi-periodic oscillations in the angle-averaged
($\Delta\theta_{cm}\simeq 90^{\circ}\pm 25^{\circ}$) excitation functions for
the ${}^{24}Mg+^{24}Mg$ elastic-inelastic scattering and $\alpha$-transfer
channels on the energy interval $E_{cm}=44.86-47.76$ MeV. The period of the
energy structures, $\simeq$0.81 MeV, is interpreted as inverse half-period
($\simeq 5\times 10^{-21}$ sec.) of coherent rotation of highly excited short-
lived ($\simeq 3.6\times 10^{-21}$ sec.) chain of a length $\simeq 30$ fm. The
rotational wave packet coherence survives (i) the energy relaxation (fully
mixing ergodic dynamics) for the strongly overlapping states with fixed total
spin and parity values, and (ii) the strong perturbation of the motion upon a
change of the total spin. The present discussion rises a number of the
questions. For example, is rotational coherence of large molecules necessarily
destroyed in the conventionally statistical limit of structureless non-
selective continuum (for fixed total spin and parity values) under the
conditions of complete intramolecular energy redistribution and vibrational
dephasing in the regime of strong ro-vibrational coupling? For the slow cross-
symmetry phase relaxation, quantum coherent superpositions of a large number
of complex configurations with, e.g., many different total angular momenta
produce image of a rotation of macroscopic object with classically fixed
(single) total angular momentum. Suppose that the quantum coherent
superpositions involving a very large number of different good quantum numbers
play a role, in a hidden form, in a formation of macroscopic world. Then why
these quantum superpositions are so stable against quick aging/decay of
ordered complex structures preventing or slowing down tendencies towards
uniform occupation of the available phase space as prescribed by the random
matrix theory? And what kind of complex macroscopic phenomena may reveal
traces of partially coherent quantum superpositions involving a huge number of
quantum-mechanically different integrals of motion behind of what is referred
to as conservation laws in classical physics employed for the description of
the macroscopic world?
###### keywords:
coherent rotation , phase relaxation , ${}^{24}Mg+^{24}Mg$ collision , nuclear
molecules , quantum macroscopic transition
Femtochemistry provides powerful tools in studying molecular structure and
dynamics in chemistry and biology [1], [2]. Nowadays, the use of ultrashort
femtosecond and attosecond laser pulses is a widely used technique for the
real-time monitoring of chemical reactions. One of the successes of
femtochemistry is its ability for direct probe of vibrational and rotational
wave packets and their coherent evolution as classical-like, essentially due
to the quantum interference, system-specific trajectories of a transitional
state of a given chemical reaction. Precondition for survival of the
characteristic time and length scales for the system-specific wave packet
dynamics is that the isolated intermediate complex is not in ergodic state.
This means that phase relaxation (dephasing) or dynamical decoherence [3],
responsible for a spread of the wave packets, is relatively slow, not faster
than the life time of the intermediate complex. This regime of incomplete
mixing is beyond applicability of random matrix theory [4] which addresses a
featureless universal behavior of complex systems in the ergodic limit of
complete lose of phase memory and uniform occupation of the available phase
space (simple pictorial explanation of the domain of applicability of random
matrix theory, for not experts in the field, is given in Sect. V of [5]).
Unfortunately, the femtochemistry experimental methods are inapplicable to
study reactions which are not initiated by the laser pulses. In particular,
time evolution of heavy ion collisions is not accessible for the direct
monitoring. Additional experimental difficulty to measure time scales of
reactions involving relatively light heavy ions is their short duration.
Namely, time resolution of order of zeptosecond (10-21 sec) would be needed
for an accurate reconstruction of heavy ion collisional dynamics. Therefore,
the only possible way to access time evolution of heavy ion reactions is
provided by the detailed energy dependence of the cross sections, i.e.
excitation functions. In this respect, the most favorable situation is the
presence of both fast (direct) and time delayed processes [6], [7], [8], [9].
Then, under the condition of a major contribution of direct processes, the
energy variations in the excitation functions are mainly given by the
interference between the collision amplitudes corresponding to the direct and
time delayed processes. Since the direct process amplitude is energy smooth
function, this allows to obtain time dependent amplitude, ${\cal P}(t,\theta)$
of the time delayed process [6], [7], [8], [9]. Clearly, the direct process
plays a role of the pump pulse switching on the clock at the initial moment of
time [8], [9]. When, for that or another reason, the interference term between
the amplitudes of the direct and time delayed processes does not contribute to
the cross section, the energy variations in the excitations functions
originate from the modulus square of the energy oscillating around zero
amplitude, $F(E,\theta)$, corresponding to the time delayed processes. In this
situation, an information on the time evolution, though not as complete as in
the presence of direct processes, is still encoded in the incident energy
dependence of the cross sections. Indeed, we represent the cross section
$\sigma(E,\theta)=|F(E,\theta)|^{2}$ as
$\displaystyle\sigma(E,\theta)\propto\int_{-\infty}^{\infty}d\tau\exp(iE\tau/\hbar)<{\cal
P}(t+\tau/2,\theta)$ $\displaystyle{\cal P}(t-\tau/2,\theta)^{\ast}>_{t}$ (1)
with
$\displaystyle<{\cal P}(t+\tau/2,\theta){\cal
P}(t-\tau/2,\theta)^{\ast}>_{t}=\int_{0}^{\infty}dt$ $\displaystyle{\cal
P}(t+\tau/2,\theta){\cal P}(t-\tau/2,\theta)^{\ast}.$ (2)
In the above relations, $E>0$ is the total excitation energy, $t$ is the time,
$\theta$ is the emission (scattering) angle and ${\cal
P}(t,\theta)\propto\int_{E_{1}}^{E_{2}}dE\exp(-iEt/\hbar)F(E,\theta)$ with
${\cal P}(t\leq 0,\theta)=0$ has a finite time resolution, $\Delta t\simeq
2\pi\hbar/(E_{2}-E_{1})$.
Suppose that a pronounced quasi-periodic component, with a period of $\delta
E$, is present in $\sigma(E,\theta)$. From model independent Eqs. (1) and (2),
this implies quasi-periodicity or recurrences, at least one recurrence, in
${\cal P}(t,\theta)$ with a period of $\simeq 2\pi\hbar/\delta E$. If, in
addition, the quasi-periodic energy components correlate for given well
resolved different exit channels, this implies that $<{\cal
P}(t+\tau/2,\theta){\cal P}(t-\tau/2,\theta)^{\ast}>_{t}$ is significantly
enhanced at approximately the same, for these different correlated channels,
moments $\tau\simeq 2M\pi\hbar/\delta E$ with $M=1,2,...$, at least with
$M=1$.
Quasi-periodic structures in the excitation functions of heavy ion collisions
often reveal presence of non-statistical effects. For example, the oscillating
energy structures in the ${}^{24}Mg+^{24}Mg$ elastic-inelastic scattering
[10], [11] were interpreted in terms of the highly excited coherently rotating
hyper-deformed intermediate complex [12]. However, the identification of the
non-statistical quasi-periodic structures from the analysis of the cross
section energy autocorrelation functions in [12] may look questionable for the
following reason. In the statistical regime of random phases of the strongly
overlapping resonance levels with different total spin values, the energy
autocorrelation functions, $C(\varepsilon)$, have a Lorentzian shape [13].
This result is obtained for very long energy interval on which excitation
functions are measured. Yet, for the finite data range analysis of the
concrete experimental data sets, one encounters deviations from the
theoretical prediction, in particular, in the form of fluctuations around the
Lorentzians [14]. Therefore, even though the oscillations in
$C(\varepsilon)$’s have approximately the same period for the different exit
channels [11], [12], which is inconsistent with the statistical
interpretation, additional supportive argumentation for the quasi-periodic
non-chaotic behavior is highly desirable. The interpretation [12] in terms of
the relatively stable coherent rotation, for the slow phase relaxation, was
based on the assumption of the highly excited ($\simeq 10-15$ MeV above Yrast
line) intermediate complex with strongly overlapping resonances. Therefore, it
should be insightful to also test this interpretation for transfer channels in
the ${}^{24}Mg+^{24}Mg$ collision. Do the energy oscillating structures, with
the same quasi-period as that for the ${}^{24}Mg(^{24}Mg,^{24}Mg)^{24}Mg$
elastic-inelastic scattering, present in the transfer channels? Affirmative
answer to this question would provide additional supportive indication that
both ${}^{24}Mg(^{24}Mg,^{24}Mg)^{24}Mg$ elastic-inelastic scattering and the
transfer reaction channels originate from the decay of the same highly excited
coherently rotating hyper-deformed intermediate complex.
The reasonings presented above has motivated us to analyze the angle-averaged,
$\Delta\theta_{cm}\simeq 90^{\circ}\pm 25^{\circ}$, excitation functions for
the ${}^{24}Mg+^{24}Mg$ elastic-inelastic scattering and $\alpha$-transfer
channels measured on the energy interval $E_{cm}=44.86-47.76$ MeV [15], which
is considerably smaller than that ($E_{cm}=42-56$ MeV) on which the data were
taken in [10], [11]. The analysis in [15] confidently confirmed the non-
statistical behavior of the energy structures in these processes. We extend
the analysis [15] concentrating on the quantitative interpretation of the
pronounced, transparently visible in Figs. 2 and 3 of [15], strongly
correlated for different channels, quasi-periodic energy structures.
Figure 1: (Color online) The normalized cross section energy autocorrelations
functions for the ${}^{24}Mg(^{24}Mg,^{20}Ne)^{28}Si$ reaction and
${}^{24}Mg(^{24}Mg,^{24}Mg)^{24}Mg$ elastic-inelastic scattering. The data
(dots) are obtained from the analysis of the summed deviation functions in
Fig. 3 of [15]. Solid lines are the fits for the relatively pure angular
resolution (though under the neglect of the interference term between the near
and far side collision amplitudes) obtained with $\hbar\omega=0.41$ MeV,
$\beta=0.02$ MeV and $\Gamma=0.2$ MeV. The fit for the poor angular resolution
(short-dashed lines) is obtained with $\hbar\omega=0.404$ MeV, $\beta=0.0205$
MeV and $\Gamma=0.17$ MeV. The Lorentzians with $\Gamma=$0.2 MeV (long-dashed
lines) are obtained for the short spin off-diagonal phase memory,
$\beta>>\Gamma,\hbar\omega$, corresponding to the limit of the random matrix
theory [4].
In Fig. 1 we present the normalized cross section energy autocorrelations
functions $C(\varepsilon)$’s constructed from the summed deviation functions
in Fig. 3 of [15]. In the upper panel of our Fig. 1, $C(\varepsilon)$ is
obtained from the three lowest transitions in the reaction
${}^{24}Mg(^{24}Mg,^{20}Ne)^{28}Si$. In the lower panel of our Fig. 1,
$C(\varepsilon)$ is obtained from addition of the three lowest well resolved
transitions in the reaction ${}^{24}Mg(^{24}Mg,^{20}Ne)^{28}Si$ and the three
lowest well resolved transitions in the ${}^{24}Mg(^{24}Mg,^{24}Mg)^{24}Mg$
elastic-inelastic scattering (see Fig. 3 in [15]). The experimental
$C(\varepsilon)$’s demonstrate oscillations with the period close to that in
Fig. 1 of [12]. Therefore, we fit the data with the formula [12]
$\displaystyle C(\varepsilon\geq 0)/C(\varepsilon=0)\simeq{\rm
Re}\\{\exp[i\pi\varepsilon/(\hbar\omega-i\beta)][1-$
$\displaystyle\exp[i\pi(\varepsilon+i\Gamma)/(\hbar\omega-i\beta)]]^{-1}\\}$
$\displaystyle/{\rm Re}\\{1/[1-\exp[-\pi\Gamma/(\hbar\omega-i\beta)]]\\}.$ (3)
Here, $\Gamma$ is total decay width of the intermediate complex, $\omega$ is a
real part of the angular velocity of the coherent rotation and $\beta/\hbar$
has a physical meaning of its imaginary part with $\beta$ being the spin off-
diagonal phase relaxation width [16]. One observes that $C(\varepsilon)$ (3)
oscillates with period $\simeq 2\hbar\omega$. The oscillations are damped for
a non-vanishing $\beta$-width. The fit (solid lines in Fig. 1) is obtained
with $\hbar\omega=0.41$ MeV, $\beta=0.02$ MeV and $\Gamma=0.2$ MeV. In the
limit of short phase memory, $\beta>>\Gamma,\hbar\omega$, $C(\varepsilon)$ has
a Lorentzian shape [13], [4] (long-dashed lines in Fig. 1 with $\Gamma=$0.2
MeV).
A few comments on the derivation and physical interpretation of Eq. (3), in
relation to the analyzed data, are in order.
1) The starting point for the derivation of Eq. (3) is a diagonal
approximation, with respect to the levels of the intermediate complex, for
$S$-matrix elements which can be justified under the condition of the fast
energy redistribution-relaxation within the states with fixed total spin and
parity values [5]. This diagonal approximation of the Bethe unitary $S$-matrix
is valid under the Bethe random sign assumption for partial width amplitudes
with given fixed total spin and parity values in a regime of strongly
overlapping resonances providing the partial decay widths are small as
comparing with the average level spacing of the highly excited intermediate
complex, see Section 56D in [17]. It is important for the present
consideration that, unlike the shell-model approach (Section IV(A2) in [4]),
the formulation [17] leads to the Bethe unitary $S$-matrix which is not
restricted to the nucleon and $\gamma$-channels only but is also applicable
for binary collisions of composite reaction partners including heavy ions.
Indeed, the physical meaningfulness of the diagonal approximation similar to
that in [17] is recognized for heavy ion induced collisions even for a
multiple nucleon cascade emission provided the reaction proceeds through an
equilibrated intermediate system [18]. We understand that the equilibration in
[18] implies (i) energy relaxation and ergodicity within the states with fixed
total angular momentum and parity values on each step of the evaporation
cascade, and (ii) the conventional assumption of a vanishing of the total spin
and parity off-diagonal correlations on each step of the evaporation cascade.
Yet, while the entrance channel orbital momentum off-diagonal interference was
neglected making directions along and opposite the incident beam
undistinguishable, the strong correlation between different orbital momenta of
the evaporated nucleons, on each step of the evaporation cascade, assumed in
[18] (Eqs. (5) and (6) in [18]) clearly contradicts to the conventional
postulate on maximal randomness of the partial width amplitudes (Eq. (2.2) in
[19]). We mention in passing that the orbital momenta correlations [18], if
exist in reality, may result in a complete uncertainty of the standard
evaluation of the spin dependence of the nuclear level densities from angular
anisotropy around 90∘ in the evaporation processes [20], [21], where such
orbital momenta correlations were ignored in a literal consistency with the
assumption on maximal randomness of the partial width amplitudes. As a next
step, we extend the argumentation [5] to justify the relations (8.1) and (8.2)
in [22] between the product of the entrance ($a$) and the exit ($b$) channel
partial width amplitudes for different $J$-values:
$\gamma_{\nu}^{Ja}\gamma_{\nu}^{Jb}=\sum_{\mu}\gamma_{\mu}^{Ia}\gamma_{\mu}^{Ib}Q_{\nu\mu}^{JI}+R_{inc},$
(4)
where
$Q_{\nu\mu}^{JI}=(1/\pi)D\beta|J-I|/[(E_{\nu}^{J}-E_{\mu}^{I}-(J-I)\hbar\omega)^{2}+\beta^{2}(J-I)^{2}]$
(5)
with $E_{\nu}^{J}$ and $E_{\mu}^{I}$ being resonance energies for the states
with total spin values $J$ and $I$, respectively. The relations (4) hold true
independent of whether the energy redistribution process is taken explicitly
into account or not [5]. We are interested in a regime of strongly overlapping
resonances, $\Gamma/D\gg 1$, where both the total decay resonance width
$\Gamma$ and the average level spacing $D$ are taken $J$-independent. Under
the condition of relatively slow spin off-diagonal phase relaxation,
$\beta\ll\Gamma$ (but still $\beta\geq D$), $R_{inc}$ in Eq. (4) can be
neglected [22] for it produces insignificant contribution into the r.h.s. of
Eq. (6). We obtain
$\displaystyle
S_{ab}^{J}(E)=W[|J-{\bar{J}}(E)|/d]^{1/2}\sum_{\mu}\gamma_{\mu}^{Ia}\gamma_{\mu}^{Ib}/[E-\hbar\omega(J-I)$
$\displaystyle-E_{\mu}^{I}+(i/2)\Gamma+i\beta|J-I|],$ (6)
where we have omitted energy smooth phase shifts which originate from the
potential mean field scattering and direct interaction taking place on the
relatively short time intervals of formation and disintegration of the
intermediate complex. In case of elastic scattering, $(a=b)$, the products
$\gamma_{\nu}^{Ja}\gamma_{\nu}^{Jb}$ and $\gamma_{\mu}^{Ia}\gamma_{\mu}^{Ib}$
in Eq. (4) and Eq. (6) should be changed to
$[(\gamma_{\nu}^{Ja})^{2}-\overline{(\gamma_{\nu}^{Ja})^{2}}^{\nu}]$ and
$[(\gamma_{\mu}^{Ia})^{2}-\overline{(\gamma_{\mu}^{Ia})^{2}}^{\mu}]$,
respectively [5]. In Eq. (6), $W[|J-{\bar{J}}(E)|/d]$ is the energy averaged
decay probability. It was taken in a bell shaped form with the maximum at
$J={\bar{J}}(E)$ and with the width $d$ of about an effective range of total
spin values coherently excited in the collision process. For the data analyzed
in this Letter this range is $\Delta J\simeq 34-38$. In the derivation of Eq.
(3), we took ${\bar{J}}(E)=I+(E-{\bar{E}})/(\hbar\omega)$, where ${\bar{E}}$
is the average energy corresponding to $J({\bar{E}})=I\gg 1$. Under the
neglect of the interference term between the near and far side collision
amplitudes due to the angle-averaging, such a choice of ${\bar{J}}(E)$ results
in $C(\varepsilon)$ (3) to be approximately independent of $d$. However, for
the other choices, e.g. ${\bar{J}}(E)=I=const$, the additional damping factor,
$\simeq W[-|\varepsilon|/(d\hbar\omega)]$, would appear in the r.h.s. of Eq.
(3).
2) The data fitted in Fig. 1 were obtained from the angle-averaged excitation
functions, $\Delta\theta_{cm}\simeq 90^{\circ}\pm 25^{\circ}$. Yet, the
formula (3) was derived for the pure angle-resolution (though under the
neglect of the interference term between the near and far side collision
amplitudes as well as between direct and time-delayed processes [12]).
Generalization for the angle-averaged excitation functions,
$\Delta\theta_{cm}\simeq 90^{\circ}\pm 25^{\circ}$, yields the additional
damping factor in the r.h.s. of Eq. (3). This factor is approximately given by
$\sin^{2}[\varepsilon\Delta\theta/(2^{3/2}\hbar\omega)]/[\varepsilon\Delta\theta/(2^{3/2}\hbar\omega)]^{2}$
with $\Delta\theta=50^{\circ}=0.88$. The corresponding fit of the data on the
lower panel of Fig. 1 is obtained with $\hbar\omega=0.404$ MeV, $\beta=0.0205$
MeV and $\Gamma=0.17$ MeV. Clearly, the additional damping factor, due to the
poor angular resolution, does not affect quasi-period of the oscillations in
$C(\varepsilon)$.
3) The time dependent intensity, $|{\cal P}(t,\theta)|^{2}$, closely relates
to the return probability for a finite time resolution, while modulus square
of Fourier component of ${\cal P}(t,\theta)$, i.e. the cross section (1), is
analogous to the low-resolution version the spectrum [23] providing the Bethe
$S$-matrix diagonal approximation is applicable (Eq. (260) and Section 56D in
[17]). In our case the finite time resolution is $\hbar/\Gamma$, and the
spectrum is not resolved due to the strong overlap of the resonance levels,
$\Gamma/D\gg 1$. Clearly, $C(\varepsilon)$ has a physical meaning of the
autocorrelation function of the low-resolution version of the spectrum [23].
In a view of the relatively short energy interval of the measurements [15],
$\Delta E_{cm}=2.9$ MeV, as comparing with $\Gamma$ the standard evaluation of
errors due to the finite data range [14] would result in the statistical
insignificance of the quasi-periodic structures in Fig. 1. However, the method
of evaluation of the statistical uncertainties [14] is not applicable in the
presence of the $S$-matrix spin off-diagonal correlations for
$\beta\ll\Gamma$. The essence of the matter can be explained as follows.
Consider a given realization of $\sigma^{I}(E)=|S_{ab}^{I}(E)|^{2}$, where
$S_{ab}^{I}(E)$ is given by Eq. (6), i.e. like in [14]. Then, in a regime of
strongly overlapping resonances, the standard evaluation of errors due to the
finite data range [14] would be certainly applicable and quasi-periodic
oscillations in the corresponding $C(\varepsilon)$, if occur at all, would not
be of the sufficient statistical significance for the short energy interval,
$\Delta E_{cm}=2.9$ MeV. Consider next
$\sigma(E)=\sum_{n=-q}^{q}\sigma^{I}(E+2n\hbar\omega)$, where $n$ and $q$ are
natural numbers and $n$ in the summation changes with the step of unity. For
$q\gg 1$, $\sigma(E)$ is close to periodic function with the period of
$2\hbar\omega$ independent of actual realization of $\sigma^{I}(E)$ the latter
being generated like in [14]. Clearly, the correspondent $C(\varepsilon\leq
r\hbar\omega)$, where $1\ll r\ll q$, will show periodic behavior even if it
constructed from $\sigma(E)$ restricted to the relatively short energy
interval $\Delta E\simeq(6-7)\hbar\omega$. For $q=1$ (the three terms in the
sum), similar to our case, $\sigma(E)$ will not be exactly periodic but will
still have the strong quasi-periodic component which will show up in the
statistically significant way in the corresponding $C(\varepsilon)$ extracted
even on the short energy interval $\Delta E\simeq(6-7)\hbar\omega$. The
numerical simulations would represent the straightforward way to demonstrate
the statistically significant registration of the quasi-periodicity for the
short energy intervals. In a view of a simplicity of the simulations these can
be performed during couple of days by a science or nowadays even, e.g.,
economy or linguistic undergraduate student. The results could be reported,
e.g., during the workshop “Information and statistics in nuclear experiment
and theory ISNET-3” at the ECT∗ on November 16-20, 2015. In addition to the
${}^{24}Mg+^{24}Mg$ collisions the simulations may also include the
${}^{12}C+^{24}Mg$ elastic and inelastic scattering which demonstrate strong
quasi-periodic structures (Fig. 1 in [7]). This is not a joke for the issue
needs clarification because when one of us (S.K.) 15 years ago discussed the
experimentally observable oscillations in Fig. 1 of [7] with A. Richter
(previous Director of the ECT∗) he was told that these oscillations are due to
the finite data range effects and the data are consistent with the random
matrix description in terms of Ericson fluctuations [4]. Then one of the
questions to be answered by the student is: How many statistically independent
generations, of the type of [14], of the excitation functions using the
standard algorithms of the random matrix theory (the standard statistical
model) [4] must be performed in order to reproduce the oscillations in Fig. 1
of [7]? Obviously the student must have an optimistic vision for the future
and be brave enough to improve the world we live in. This is because the
opposition is too strong, crossed all possible red lines and, therefore, has
nothing to lose [5]. For we are speaking about those who demonstrated the
inquisitory attitude and are killers of the new ideas which does not leave
much choice at this stage. We do not mean here and put aside the Australian
way proved to bring this kind of attitude up to the national operational
policy which, bedazzled by the ideas of universality of the random matrix
theory, is spiritually oriented on the future with no distinguishability
between living human beings and corpses [5].
The work under the extended version of this Letter, where the derivation of
the cross section energy autocorrelation (3) and some other relevant issues
will be addressed, is under way. This will include a discussion of adequacy of
the conventional interpretation of the non-statistical structures in terms of
isolated resonances. For example, what is angular dependence of
$C(\varepsilon)$, in particular, for the elastic
${}^{24}Mg(^{24}Mg,^{24}Mg)^{24}Mg$ scattering in the presence of dominant
potential scattering contribution ($\theta\leq 70^{\circ}$), predicted within
our treatment in terms of the coherent rotation of the intermediate complex
with strongly overlapping resonances as compared with that resulted from the
interpretation in terms of isolated resonances as identified in [24]?
Stability of the coherent rotation may be appreciated by noticing that the
spreading of the angular orientation during one revolution due to the finite
$\beta$-width is $2\pi\beta/\hbar\omega\simeq 15^{\circ}$. Therefore, it takes
about 10 complete revolutions, i.e. $\simeq 10^{-19}$ sec., for the complex to
loose the coherent nature of its rotation. Energy relaxation time scale for
excited nuclei is $\simeq\hbar/\Gamma_{spr}\simeq 10^{-22}$ sec., where
$\Gamma_{spr}\simeq 5$ MeV is the spreading width [4]. We observe that the
coherent rotation persists for about three orders of magnitude longer than it
takes to complete a process of energy equilibration (vibrational relaxation by
the molecular physics terminology).
The strong channel correlation between individual well resolved channels in
both the ${}^{24}Mg(^{24}Mg,^{24}Mg)^{24}Mg$ elastic-inelastic scattering and
${}^{24}Mg$ $(^{24}Mg,^{20}Ne)^{28}Si$ transfer reaction, as well as between
the channels for these two processes, is transparently visible in Fig. 2 of
[15]. Therefore, it is obvious that strongly correlated regular oscillations
with the approximately channel-independent quasi-period, $\simeq 0.81$ MeV,
are present in all the six channels. This rules out statistical origin of
these strongly correlated quasi-periodic structures and, therefore, their
interpretation in terms of Ericson fluctuations [13]. This is because Ericson
fluctuations, which is a particular case of the Bethe statistical fluctuations
(to be reported in the extended version of this Letter), produce uncorrelated,
for different channels, irregular energy structures. As a result, the
deviations from the Lorentzian must be uncorrelated for different channels
instead of demonstrating regular oscillations with the approximately channel-
independent quasi-period for all the six channels. The fact that the present
analysis of the data on the energy interval $E_{cm}=44.86-47.76$ MeV revealed
a value of the quasi-period, $\simeq 0.81$ MeV, very close to that obtained
for the much longer energy interval $E_{cm}=42-56$ MeV [12], supports our
interpretation. A possibility of further test of our interpretation, e.g. for
the ${}^{24}Mg(^{24}Mg,^{24}Mg)^{24}Mg$ elastic scattering, is suggested by
the expected dominance of direct processes (potential scattering) for
$\theta\leq 70^{\circ}$ (see Figs. 6 and 7 in [24]). Then the energy
variations in the excitation functions originate mainly from the interference
between the collision amplitudes corresponding to the direct and time delayed
processes. As a result, the characteristic quasi-periods of leading harmonics
in $\sigma(E,\theta)$ become strongly $\theta$-dependent [7], [8], [9], which
can be tested in experiments with pure angular resolution. We mean the
numerical experiments against the previously measured relevant available data
since the experimental excavation into the matter as well as the understanding
of the underlying phenomena are no longer among interests of the nuclear
physics community.
The numerical experiments are useful to illustrate simple algorithms of
transforming/processing the complex quantum information. The initial wave
packet is given by a coherent superposition of partial waves with different
orbital momenta describing the colliding objects in their ground states. In a
process of the collision the initial quantum superpositions are transformed
into the coherent superpositions of a large number of strongly overlapping
resonance configurations of the intermediate complex with different total spin
values. The spin off-diagonal correlations produce the coherent rotation of
the intermediate complex with a well defined angular velocity. Yet, the
intermediate complex represents a quantum superposition of the two objects
simultaneously rotating in opposite directions. For these coherently rotating
in opposite directions alternatives, each of them individually is itself
result of the spin off-diagonal quantum interference, produce interference
fringes [9]. For small $\beta$-width, the rotating wave packets slowly spread
and interference fringes, related to the spin off-diagonal correlations,
disappear [25]. In this way the complex collision experiments provide a
physical visualization of executing the quantum computing algorithms for a
simple description of the quantum-macroscopic transition, where the
convergence to the macroscopic-like dynamics (rotation) on the final stage of
the “calculations” is due to partial dynamical decoherence but still images
essentially from the quantum interference. Each step of the “quantum
computing”, presumably run by nature, can be easily mapped on the fundamental
double slit experiment (since the many slit experiment is a combination of the
double slit ones). The double slit experiment “has been designed to contain
all of the mystery of quantum mechanics …” [26] and, Feynman continues, “Any
other situation in quantum mechanics, it turns out, can always be explained by
saying, ’You remember the case of the experiment with the two holes? It’s the
same thing’.” [26]. The above simple picture of transformation of the complex
quantum information is noticeably supported by the state of art numerical
simulations for, e.g., H+D${}_{2}\to$HD+D molecular reaction (Fig. 1 in [27])
even though the calculations [28] rules out a picture of the isolated
resonances of the intermediate complex (Fig. 1 in [28]). We mention in passing
that the described above algorithms are incompatible with basic ideas of the
random matrix theory as applied to classically chaotic systems [4]. Therefore
it is inappropriate to rely on mentality of indistinguishability and selective
blindness of the random matrix theory [5], which unfortunately found its way
beyond the mere academic and education activities, as a point of reference in
discussing the above mentioned system specific algorithms.
Angular velocity of coherent rotation, $\omega$, can not be defined within the
states with fixed $J$-value. This transparently follows from the conjecture
[5] which relates $Q_{\nu\mu}^{JI}$ (5) to the spin off-diagonal correlations
between squares of the individual resonance wave functions of the intermediate
complex. Indeed, it is obvious that distribution of the spacial density,
including the radial extension, is different for the state with fixed
$J$-value from that for the coherently rotating intermediate complex.
Therefore, energy of the coherent rotation as well as the associated moment of
inertia and, thus, the characteristic length scale of the coherently rotating
intermediate complex are also undefinable within the states with fixed
$J$-value.
Let us assume that the coherent rotation may be considered as the macroscopic
motion. Then, an assumption of the spherical intermediate complex having rigid
body moment of inertia is ruled out for, in this case, $\hbar\omega\simeq 4$
MeV, resulting in the rotational energy ($\simeq 70$ Mev) which exceeds the
total excitation energy ($\simeq 60-63$ MeV). Instead, the small value of
$\hbar\omega\simeq 0.4$ MeV, corresponding to the period of the rotation
$\simeq 10^{-20}$ sec, indicates an anomalously strong deformation of the
coherently rotating complex with $J\simeq 34-38$. We calculate moment of
inertia of this coherently rotating complex, ${\cal J}_{coh}$, and find that
it corresponds to a chain-state of the length $\simeq 30$ fm. This is close to
the length $\simeq 24$ fm of the chain of four touching carbon nuclei,
calculated with $r_{0}=1.3$ fm. We evaluate deformation energy of the chain as
a sum of Coulomb energy of the two touched ${}^{24}Mg$ nuclei and double
Coulomb energy of the two touched ${}^{12}C$ nuclei. For $r_{0}$=1.3 fm, this
gives $\simeq$38 MeV for the deformation energy. The energy of the coherent
rotation, for $J\simeq 34-38$, is $\simeq 7$ MeV. Therefore, the energy of the
intrinsic excitation (heat) is estimated $\geq 15$ MeV, i.e. our “carbon
sausages” are really hot.
Let us evaluate average level spacing ($D$) of the intermediate complex with
$J\simeq 34-38$. We use the standard statistical model formula (the Weisskopf
evaluation [29], [30], [31]) for the total decay width with fixed
$(J,\pi)$-values: $\Gamma=(1/2\pi)Dn{\bar{T}}$. Here, $n$ is a number of the
open channels and ${\bar{T}}$ is the averaged over the channels transmission
coefficient. The total number of channels, specified by the channel spins,
orbital momenta and internal quantum numbers of the fragments, for the
${}^{24}Mg+^{24}Mg$ elastic and inelastic scattering, up to excitation of the
(6+,4+) states of the ${}^{24}Mg$ collision fragments, is $57$. We estimate
number of the $\alpha$-transfer and two-$\alpha$-transfer channels and the
proton and $\alpha$-emission channels to be $\simeq 200$, where we have taken
into account all possible the channel spins and orbital momenta for given
internal quantum numbers of the fragments. Then our estimate is $n\simeq 250$.
Since the heavy ion collisions under consideration are characterized by a
closeness to the potential barrier (centrifugal plus Coulomb barriers) or even
are sub-barrier, and the proton and $\alpha$-particle evaporation is mostly
sub-barrier, we estimate ${\bar{T}}\simeq 0.5$. Then, for $\Gamma=0.2$ MeV, we
obtain $D\simeq$10 keV. This indicates the relatively weak coupling to the
continuum, $\Gamma/(nD)\simeq 0.08$, in the regime of strongly overlapping
resonances, $\Gamma/D\simeq 20$. Therefore the Bethe diagonal approximation
(Section 56D in [17]) obtained from the Bethe unitary $S$-matrix for binary
reactions involving complex collision partners is applicable (to be presented
in the extended version of this Letter).
Note that the standard statistical model evaluation [17], [31], [32], for the
intrinsic excitation of 15 MeV of the intermediate nucleus, yields
$D\simeq$0.01-0.1 keV. This is 2-3 orders of magnitude smaller than the
average level spacing evaluated above. We interpret this reduction of $D$ as
an indication of a sizable admixture of clustering configurations in the wave
function of the highly excited strongly deformed intermediate complex with
large spin values $J\simeq 34-38$. These clustering configurations should not
be considered as vibrational modes of the linear chain of four carbon nuclei
each being in its ground state. Indeed it seems hardly probable that such a
configuration would be stable with respect to decay. Therefore, it is perhaps
more realistic to consider the hyper-deformed highly excited states as the
approximately linear chain of four relatively heavy clusters in combination
with some “valent” nucleons and $\alpha$-particles.
Up to now the idea of rotational coherence of ergodic, with respect to the
energy and phase relaxation within the states with fixed $J$-values, nuclear
molecules in continuum has been employed to extract spectroscopic information
($\hbar\omega,{\cal J}_{coh}$) from dynamics for a few colliding systems,
${\sl e.g.}$, ${}^{58}Ni+^{58}Ni$ [33], ${}^{58}Ni+^{62}Ni$ [33], [34],
${}^{46}Ti+^{58}Ni$ [34], ${}^{12}C+^{24}Mg$ [6], [7], [8], [9],
${}^{24}Mg+^{24}Mg$ and ${}^{28}Si+^{28}Si$ [12], ${}^{24}Mg+^{28}Si$ [35].
Therefore, it may be of interest to apply this method for analysis of a large
number of available data sets for many heavy ion colliding systems. On the
other hand, the question is whether the stable coherent rotation survives
violation of rotational symmetry and other external perturbations of the
Hamiltonian? The conjecture [5] offers a possible specification of this
question in terms of the correlation properties of eigenstates of the
intermediate complex.
Experimental manifestation of rotational coherent motion in complex molecules
was initially met with scepticism “because of the general belief that Coriolis
interactions, anharmonicity and other interactions would destroy the
coherence.” (see subsection “Changing a Dogma: Development of RCS” in [2]).
The present discussion leads us to ask: Is rotational coherence of complex
molecules necessarily destroyed in the conventionally statistical ergodic
limit of structureless (non-selective) continuum under the conditions of
complete intramolecular energy redistribution (thermalization) and vibrational
dephasing in the regime of strong ro-vibrational coupling? Though challenging
experimentally, the subject is of interest to search for fingerprints of
transformation of the quantum coherent wave packet dynamics into macroscopic
motion. More specifically, in the quantum regime, the spreading of the wave
packets is determined by $\beta$-width, which has essentially quantum origin
of the spin off-diagonal correlations [5]. Then what is a classical analog of
$\beta$-width, for classically chaotic systems, when macroscopic non-linear
dynamics takes over from the coherent rotation of the spreading quantum wave
packets? Addressing the questions above should include a comparison of the
nature of the dephasing analyzed in [1] with that discussed in this Letter
(see also [5] and references therein).
## Acknowledgments
We thank Sergei Maydanyuk for useful discussions. Writing this Letter started
when one of us (S.K.) was supported by Chinese Academy of Sciences Visiting
Professorship for Senior International Scientists (2012-2013).
## References
* [1] A. Zewail, Femtochemistry, Vols. I and II, World Scientific, Singapore, 1994.
* [2] A. Zewail, www.nobelprize.org/nobel-prizes/chemistry/laureates/1999/zewail-lecture.pdf, and references therein.
* [3] G. Casati, B.V. Chirikov, Phys. Rev. Lett. 75 (1995) 350.
* [4] G.E. Mitchell, A. Richter, H.A. Weidenmüller, Rev. Mod. Phys. 82 (2010) 2845, and references therein.
* [5] S. Kun, Y. Li, H. Zhao, M.R. Huang, e-print arXiv: quant-ph/1307.4490.
* [6] S.Yu. Kun, A.V. Vagov, O.K. Vorov, Phys. Rev. C 59 (1999) R585.
* [7] S.Yu. Kun, A.V. Vagov, W. Greiner, Phys. Rev. C 63 (2001) 014608.
* [8] S.Yu. Kun, A.V. Vagov, L.T. Chadderton, W. Greiner, Int. J. Mod. Phys. E 11 (2002) 273.
* [9] L. Benet, S.Yu. Kun, Wang Qi, Phys. Rev. C 73 (2006) 064602, e-print arXiv: quant-ph/0503046.
* [10] R.W. Zurmühle, P. Kutt, R.R. Betts, S. Saini, F. Haas, Ole Hansen, Phys. Lett. B 129 (1983) 384.
* [11] A. Sarma, R. Singh, Z. Phys. A 329 (1988) 195.
* [12] S.Yu. Kun, B.A. Robson, A.V. Vagov, Phys. Rev. Lett. 83 (1999) 504.
* [13] T. Ericson, T. Mayer-Kuckuk, Ann. Rev. Nucl. Sci. 16 (1966) 183.
* [14] P.G. Dallimore, I. Hall, Nucl. Phys. 88 (1966) 193.
* [15] S. Saini, R.R. Betts, R.W. Zurmühle, P.H. Kutt, B.K. Dichter, Phys. Lett. B 185 (1987) 316.
* [16] S.Yu. Kun, Z. Phys. A 357 (1997) 271.
* [17] H.A. Bethe, Rev. Mod. Phys. 9 (1937) 69.
* [18] S.E. Koonin, W. Bauer, A. Schäfer, Phys. Rev. Lett. 62 (1989) 1247.
* [19] H. Feshbach, A. Kerman, S.E. Koonin, Ann. Phys. 125 (1980) 429.
* [20] V.M. Strutinsky, in Proceeding of Conference on Nuclear Reactions at Low and Medium Energies, Moscow 1957, p. 522, Academy of Sciences USSR Publishing, 1958.
* [21] T. Ericson, V.M. Strutinsky, Nucl. Phys. 8 (1958) 284; 9 (1958/1959) 689.
* [22] S.Yu. Kun, Z. Phys. A 357 (1997) 255.
* [23] E.J. Heller, Phys. Rev. Lett. 53 (1984) 1515.
* [24] A.H. Wuosmaa, R.W. Zurmühle, P. Kutt, S.F. Pate, S. Saini, M.L. Halbert, D.C. Hensley, Phys. Rev. C 41 (1990) 2666.
* [25] L. Benet, L.T. Chadderton, S.Yu. Kun, Wang Qi, Phys. Rev. A 75 (2007) 062110; e-print arXiv: quant-ph/0610091.
* [26] R. Feynman, The Character of Physical Law, based on the Messenger Lectures given by R. Feynman at Cornell University in 1964, Lecture 6: Probability and Uncertainty - the Quantum Mechanical view of Nature, first published by BBC, 1965.
* [27] P.D.D. Monks, J.N.L. Connor, S.C. Althorpe, J. Phys. Chem. A 110 (2006) 741.
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* [30] J.M. Blatt, V.F. Weisskopf, Theoretical Nuclear Physics, Wiley, 1952.
* [31] L. Landau, Ya. Smorodinsky, Lectures on Theory of Atomic Nucleus, Moscow, Technical-Theoretical Publishing, 1955.
* [32] A. Bohr, B. Mottelson, Nuclear Structure, Vol. 1, Benjamin, New York, 1969.
* [33] L. Vannucci, U. Abbondanno, M. Bettiolo, M. Bruno, N. Cindro, M. D’Agostino, P.M. Milazzo, R.A. Ricci, T. Ritz, W. Scheid, G. Vannini, Z. Phys. A 355 (1996) 41.
* [34] S.Yu. Kun, U. Abbondanno, M. Bruno, N. Cindro, M. D’Agostino, P.M. Milazzo, R.A. Ricci, T. Ritz, B.A. Robson, W. Scheid, A.V. Vagov, G. Vannini, L. Vannucci, Z. Phys. A 359 (1997) 145.
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|
arxiv-papers
| 2013-09-25T02:49:05 |
2024-09-04T02:49:51.478876
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "M. H. Zhao, S. Kun, O. Merlo, M. R. Huang, Y. Li, J. S. Wang",
"submitter": "Meirong Huang",
"url": "https://arxiv.org/abs/1309.6386"
}
|
1309.6494
|
# BENCHMARKING HEADTAIL WITH ELECTRON CLOUD
INSTABILITIES OBSERVED IN THE LHC
H. Bartosik W. Höfle G. Iadarola Y. Papaphilippou G. Rumolo
CERN Geneva Switzerland
###### Abstract
After a successful scrubbing run in the beginning of 2011, the LHC can be
presently operated with high intensity proton beams with 50 ns bunch spacing.
However, strong electron cloud effects were observed during machine studies
with the nominal beam with 25 ns bunch spacing. In particular, fast transverse
instabilities were observed when attempting to inject trains of 48 bunches
into the LHC for the first time. An analysis of the turn-by-turn bunch-by-
bunch data from the transverse damper pick-ups during these injection studies
is presented, showing a clear signature of the electron cloud effect. These
experimental observations are reproduced using numerical simulations: the
electron distribution before each bunch passage is generated with PyECLOUD and
used as input for a set of HEADTAIL simulations. This paper describes the
simulation method as well as the sensitivity of the results to the initial
conditions for the electron build-up. The potential of this type of
simulations and their clear limitations on the other hand are discussed.
## 1 INTRODUCTION
At the early phase of the 2011 LHC run, seven days were devoted to scrubbing
using the 50 ns bunch spacing beam with gradually increasing number of bunches
circulating in the machine. This allowed to sufficiently condition the inner
surface of the LHC beam screens and vacuum chambers for running with the same
beam for physics production in routine operation throughout 2011. First
attempts to inject a beam with the nominal 25 ns bunch spacing were performed
at the end of June 2011 with bunch trains of 24 bunches. Pressure rise and
increased heat load in the arcs were observed. The studies continued with
injections of 48 bunches in August, where the beams became transversely
unstable after about 1000 turns with the transverse damper switched on and
after about 500 turns without transverse damper, for two injections
respectively. As will be discussed in more detail below, an analysis of the
turn-by-turn bunch-by-bunch data from the transverse damper pick-ups [1]
points to the observation of coherent electron cloud instabilities. During
these injection tests the chromaticity was set to about $Q^{\prime}\approx 2$
in both planes, as usually used on the LHC flat bottom in routine operation.
Only after increasing the chromaticity to about $Q^{\prime}\approx 15$ in both
planes during further studies in October, it was possible to perform nominal
injections of 288 bunches with 25 ns bunch spacing from the SPS. The high
chromaticity suppressed the fast instabilities observed before. However, the
beam suffered from slow losses and transverse emittance blow-up along the
bunch train and, as before, pressure rise and increased heat load were
measured in the cold arcs. The conditioning of the LHC beam screens due to the
beam based electron bombardment was demonstrated by estimating the secondary
electron yield (SEY) from a comparison of the measured heat load data in the
arcs with PyECLOUD simulations [2, 3]. Figure 1 shows the evolution of the
obtained maximum SEY ($\delta_{\text{max}}$) as a function of beam time in the
LHC together with the total intensity for both beams [4]. Since the heat load
data can be read only per half cell in the arcs, the SEY can be estimated only
for both beams at the same time (grey markers) unless there is only one beam
present in the machine (red markers for beam 2). A clear conditioning effect
from $\delta_{\text{max}}\\!=\\!2.1$ at the end of June when running with 50
ns bunch spacing to $\delta_{\text{max}}\\!=\\!1.52$ after approximately 50 h
beam time with 25 ns bunch spacing is observed. The analysis and simulation
studies presented in this paper will concentrate on the aforementioned
injection tests with 48 bunches in August, where $\delta_{max}$ was estimated
to be still around 2.1 (cf. Fig. 1).
Figure 1: Evolution of $\delta_{max}$ in the LHC arcs as obtained by
reproducing the measured heat load on the LHC beam screens using PyECLOUD [4].
## 2 Analysis of turn-by-turn data
Figure 2: Transverse oscillations measured with the pick-ups of the transverse
feedback for 73 turns close before the beam dump, for the first case with
transverse damper (left) and the second case without transverse damper (right)
in the horizontal (top) and vertical plane (bottom).
As mentioned before, a first attempt to inject bunch trains of 48 bunches
spaced by 25 ns into the LHC with the an intensity of about
$1.0\\!\times\\!10^{11}$ p/b was made for beam 2 during a machine development
session on August 26, 2011. Shorter bunch trains of 12 and 24 bunches with the
same bunch spacing had been successfully injected earlier in the session. It
was planned to establish the injection of longer bunch trains and complete the
transverse damper set-up for the nominal bunch spacing of 25 ns. A first
injection of 48 bunches with the transverse damper switched on was dumped
after 1000 turns due to a beam excursion interlock (referred to as dump 1).
Another injection of 48 bunches without transverse damper was aborted already
after around 500 turns due to a beam loss interlock (referred to as dump 2).
In both cases, chromaticity was set to the values usually used during
operation ($Q^{\prime}\approx 2$).
Data of the damper pick-ups for about 73 turns are stored for both beams in
the post mortem system at a beam abort. Each data set (channel) represents one
of the allowed 3564 bunch positions per ring. Figure 2 shows the turn-by-turn
oscillation of each of the 48 bunches around the closed orbit in both planes
for the 2 injections just before the respective beam dump. The oscillation
amplitude is very small for the first 25 bunches, especially when the damper
is switched on. In this case the last bunches of the bunch train reach peak
values of around 1 mm in the horizontal plane and up to 6 mm in the vertical
plane, thus the instability is mainly observed in the vertical plane. Without
damper, the last bunches of the train exhibit oscillations of up to 3 mm
amplitude in both planes. These observations are compatible with a coupled
bunch instability (slightly stronger in the horizontal plane than in the
vertical plane), which can be suppressed by the damper, and a high frequency
instability with a broad spectrum mainly in the vertical plane. The frequency
spectrum without damper is dominated by coupled bunch modes up to 1 MHz in the
horizontal plane and up to 2 MHz in the vertical plane. With damper on,
instabilities in both planes are damped up to a frequency of about 15 MHz [1].
This can be explained by a “single bunch” instability in the vertical plane
and a coupled bunch instability (mainly) in the horizontal plane, as has been
observed in the past already in the SPS [5]. The beam observations during the
injection tests in the LHC are used to benchmark numerical simulations as
described in the following sections.
## 3 Simulation procedure
Figure 3: Electron density per unit length along the bunch train as obtained
from a PyECLOUD simulation. A snapshot of the electron distribution before
each bunch entrance as indicated by the red dots is used as input for
individual HEADTAIL simulations.
The above described experimental observations were studied by combining
PyECLOUD [2] and HEADTAIL [6] simulations. As first step the electron cloud
build-up in the LHC dipoles is simulated with PyECLOUD for an r.m.s. bunch
length of 12 cm and equal bunch intensities of $N=1.0\times 10^{11}$ p/b along
the train111It was not possible to retrieve the actual bunch-by-bunch
intensity variation from the data logging during the injection tests due to
the short time of circulating beam before the dump.. Figure 3 shows the
obtained electron density per unit length along the bunch train for a
secondary electron yield of $\delta_{\text{max}}\\!=\\!2.1$. The simulation
was initialized with a uniform distribution using the same number of primary
electrons as in the simulations for determining the evolution of the SEY
discussed above [3]. The importance of the number of primary electrons and
their spatial distribution for the onset of the e-cloud instability will be
discussed in more detail in the next section. The red markers indicate the
time steps in the simulation just before the bunch passages. The electron
distributions sampled at these points are used as input for a set of 48
HEADTAIL simulations, one for each bunch.
Figure 4: Electron distribution before the passage of the last bunch of the
train as obtained from the PyECLOUD simulation. Note that the cross section of
the LHC beam screen is approximated as ellipse.
Figure 4 shows for example the electron distribution before the entrance of
the last bunch of the train. Note the accumulation of electrons in two
stripes, as typically obtained for the build-up with 25 ns bunch spacing in a
geometry as the LHC beam screen when the electrons move in a strong dipolar
magnetic field. In HEADTAIL the electron cloud is represented by thin slices
lumped at several accelerator sections and the electron motion is frozen in
the horizontal plane in order to account for the effect of the dipole magnetic
field. At each electron cloud section the electron proton interaction is
computed consecutively for longitudinal bunch slices. After a complete bunch
passage the electron cloud is reset to the initial distribution for the next
interaction section. Figure 5 shows the evolution of the transverse emittance
for a few selected bunches in the middle of the bunch train obtained with
HEADTAIL for the simulation of 500 turns at the LHC injection plateau using
the initial electron distribution as described above.
Figure 5: Evolution of the horizontal (top) and vertical (bottom) emittance
along 500 turns for a few selected bunches in the center of the bunch train.
No instability or coherent excitation of the horizontal bunch motion is
observed, as was expected since the electron cloud is located in dipole
regions. On the other hand the electron cloud drives a single bunch
instability in the vertical plane. This instability can be observed in the
form of an exponential emittance growth, as in the case studied here for all
bunches after bunch number 25. Figure 6 shows the bunch-by-bunch oscillations
in the vertical plane during the instability (from turn 50 to turn 120) as
obtained from HEADTAIL.
Figure 6: Bunch-by-bunch oscillations in the vertical plane as obtained with
HEADTAIL. For better comparison, only 70 turns during the instability (from
turn 50 to turn 120) are shown.
The simulation is in good agreement with the experimental observations during
the first injection of 48 bunches in August with the transverse damper on,
where the second half of the bunch train is unstable mainly in the vertical
plane (cf. Fig. 2), as the damper is suppressing horizontal coupled bunch
instabilities. Horizontal instabilities are not observed in the simulation,
since the coupling between bunches is not taken into account but each bunch is
treated by an independent HEADTAIL simulation. To include this coupling would
require a combined simulation of PyECLOUD and HEADTAIL, which is envisaged to
be implemented at a later stage. This will allow to reproduce better the
instability observed during the injection without transverse damper.
## 4 Sensitivity studies
It was shown in the previous section that the instability observed at LHC
injection can be reproduced with the presented simulation procedure. In the
following the sensitivity of the simulation result to the number and
distribution of the primary electrons will be studied. Furthermore, it is
observed that the electron cloud is formed mainly in two symmetric vertical
“stripes” around the beam center (cf. Fig. 4) and only a small number of
electrons is concentrated in the central part of the chamber. Thus it is
interesting to assess which part of the electron cloud distribution is mainly
responsible for driving the beam unstable.
### 4.1 Electron density
Figure 7 shows an example for a bunch in the middle of the train, which
encounters an electron density which is just above the instability threshold.
The top graph shows a histogram of the horizontal electron distribution before
the bunch arrival as obtained from the PyECLOUD simulation and used as input
for the HEADTAIL simulation. Note again
Figure 7: From top to bottom: Histogram of the horizontal electron
distribution; transverse electron distribution divided into colored regions;
evolution of the vertical emittance for the interaction with the electrons of
the respective colored area; evolution of the vertical emittance for different
horizontal cuts of the electron distribution.
that most of the electrons are concentrated in the two stripes far away from
the beam (the $3\,\sigma$ beam envelope covers roughly the two central bins).
The graph below shows for the same case the actual positions of the macro-
electrons in the transverse plane. The macro-particle distribution is divided
here into five regions as indicated by the color code (corresponding to the
bins in the histogram). In the following, the contribution from these five
regions to the vertical instability is studied. First, each of these regions
is considered independently and used as input for individual HEADTAIL
simulations. The evolution of the vertical emittance for these five cases is
shown for 500 turns in the third graph of Fig. 7. In the case studied here,
only the central part of the electron distribution is able to drive the beam
unstable. The bottom graph shows the vertical emittance evolution in case all
electrons enclosed by the respective colored areas have been removed, i.e.
cutting the electron distribution at the inner borders of the colored regions.
The instability appears only when the electrons in the central part of the
distribution are taken into account. It follows that for the typical electron
cloud distributions encountered in the simulations of the LHC dipole regions,
it is mainly the central electron density which determines the onset of the
instability. Therefore the instability threshold can be inferred roughly from
the central electron density. Only in cases where the central electron density
is _very_ small compared to the density in the stripes this approximation may
not hold.
In the following, the central density of the electron distribution is thus
used to determine the instability threshold. It should be noted that the
instability thresholds found in this manner are consistent with thresholds
previously found in HEADTAIL simulations assuming a uniform electron
distribution before the bunch passage, which is around $1\times
10^{12}\,e^{-}/$m3 [7].
### 4.2 Dependence on number of primary electrons
One of the main uncertainties of the electron cloud build-up simulations for
the LHC injection energy is the number of primary (seed) electrons. Since the
synchrotron radiation at 450 GeV is inefficient to generate photo electrons,
it is assumed that the primary electrons are created by rest gas ionization.
Therefore the number of primary electrons depends on the pressure. However the
static pressure without beam is much lower compared to the pressure levels
measured after beam injection. Figure 8 shows the central density along the
bunch train of 48 bunches for $\delta_{\text{max}}=2.1$ and two different
cases: for a small number of seed electrons equivalent to the static pressure
in the LHC cold sections (top) and for a large number of seed electrons
(bottom). The dashed red line shows the central electron density at the
instability threshold. As expected, it takes more bunch passages to build up
the electron cloud beyond the instability threshold for a smaller number of
primary electrons. The onset of the instability along the bunch train depends
thus on the pressure level assumed in the build-up simulation. For
completeness it should be emphasized that the saturation level of the central
density is very similar in both cases (as it depends mostly on
$\delta_{\text{max}}$).
Figure 8: Comparison of the central electron density along the bunch train for
a number of seeding electrons comparable to the static pressure in the LHC
dipole sections (top) and for a larger number of seeding electrons (bottom).
### 4.3 Dependence on primary electron distribution
Figure 9: Example for an electron cloud build-up simulation assuming a uniform
transverse distribution of the primary electrons all across the vacuum
chamber: central density along the bunch train (left), horizontal distribution
for bunch 33 (middle) and for the last bunch (right).
Figure 10: Example for an electron cloud build-up simulation assuming a
Gaussian distribution of the primary electrons on top of a 10% uniform
background all across the vacuum chamber: central density along the bunch
train (left), horizontal distribution for bunch 33 (middle) and for the last
bunch (right).
In addition to the number of primary electrons, also their distribution is not
well known in the case of the LHC at injection energy. In order to reproduce
the experimentally observed vertical stripes of the electron distribution in
saturation, primary electrons have to be assumed across the entire cross
section of the vacuum chamber (since the electrons are bound by the vertical
magnetic field lines). Therefore a uniform distribution of the primary
electrons can be optionally used in PyECLOUD simulations. Figure 9 shows an
example for a uniform distribution of the seed electrons together with the
histogram of the horizontal electron distribution for two selected bunches (a
bunch in the second half of the train where the saturation level is reached,
and the last bunch). The central density is very similar in the two cases, but
the stripes build-up further outside towards the end of the bunch train.
Considering that the generation of the primary electrons is caused by rest gas
ionization, it can be argued that the primary electrons should follow a
Gaussian distribution similar to the proton beam. In this case the formation
of the stripes can be achieved by adding about 10% of the total number of
electrons in the form of a uniform background all across the vacuum chamber.
Figure 10 shows the evolution of the central electron density along the bunch
train assuming such a distribution of seed electrons together with the
horizontal electron distribution for selected bunches as in Fig. 9. Here a
smaller number of seed electrons was chosen in order to obtain the instability
onset roughly at the same bunch along the bunch train as for the case of the
uniform distribution of seed electrons. However, the evolution of the electron
distribution along the bunch train is quite different for the two cases. In
comparison to the case of the uniform distribution, the central density
reaches much higher peak values in the saturated part of the train and it
takes longer to develop the stripes in the outer part of the vacuum chamber
assuming the Gaussian distribution. This can be understood intuitively, since
the center of the chamber is seeded with a larger number of electrons while
the outer regions of the vacuum chamber are seeded with a comparably smaller
number of electrons.
## 5 Compatible parameter space
In fact the number and distribution of the primary electrons are among the
biggest uncertainties in the present understanding of the electron cloud
build-up in the LHC dipoles at injection energy. On the other hand, the
estimation of $\delta_{\text{max}}$ from the reproduction of the measured heat
load in the LHC beam screens with PyECLOUD [4] is not so sensitive to the
assumptions on the seed electrons if the bunch train is long and the electron
saturation level is reached for many bunches. The estimation of
$\delta_{\text{max}}$ for the case studied here was not done for the same time
as the observations of the injections of the 48 bunches and therefore might be
slightly smaller than $\delta_{\text{max}}=2.1$. In the following, the
available parameter space is thus scanned for compatible solutions reproducing
the observed instability at LHC injection in August 2011. Figure 11 shows
compatible solutions, i.e. cases where the onset of the instability is close
to bunch 25 (between bunch 23 and bunch 27), for a range of possible values
for $\delta_{\text{max}}$ and for the number of primary electrons per bunch
passage. The two colored regions correspond to the cases of a purely uniform
distribution of the primary electrons and a Gaussian distribution with a 10%
uniform background, respectively. As expected, higher pressure levels (or
equivalently a larger number of primary electrons) and a larger secondary
electron yield is needed to reproduce the observations in the case of a
uniform electron distribution, due to the smaller central electron density
compared to the case with a Gaussian electron distribution. Note that the
range of the number of seed electrons per bunch passage explored in this study
corresponds to room-temperature equivalent pressures between 10 nTorr and 200
nTorr (the static pressure in the LHC cold sections is around 32 nTorr room-
temperature equivalent). Further studies will be needed in the future in order
to better understand the mechanism of the primary electron production and
their distribution at LHC injection energy.
Figure 11: Range of parameters reproducing the observed instability in the
LHC, as indicated by the colored areas assuming a uniform primary electron
distribution (blue) and the Gaussian distribution on top of a 10% background
(green).
## 6 CONCLUSIONS AND OUTLOOK
Fast instabilities were observed during the first two attempts of injecting
bunch trains of 48 bunches with the nominal 25 ns bunch spacing in August
2011. The onset of vertical single bunch instability can be reproduced in good
agreement by electron cloud simulations using a combination of PyECLOUD for
the build-up and HEADTAIL for the beam dynamics part. The simulations are
based on values of $\delta_{\text{max}}$ as estimated from the measured heat
load data in the LHC using PyECLOUD. The presented studies based on a
combination of PyECLOUD with the HEADTAIL code can therefore be considered as
consistency check of the current model of the electron-cloud effects on the
LHC flat bottom. The simulations show a vertical instability, fast emittance
growth and strong losses, similar to observations in the LHC. The instability
onset depends strongly on the central electron density seen by the beam, which
itself depends on the initial conditions assumed for the build-up, namely the
number of primary electrons and their distribution. Future studies should aim
at improving the understanding of the mechanism responsible for the generation
of the primary electrons. This will become important also for the estimation
of $\delta_{\text{max}}$ from the measured heat load close to the threshold of
the electron cloud build-up, despite the fact that the beam is not expected to
suffer from the instability any more. Apart from this, it would be interesting
to include the coupling between the bunches due to the electron cloud in the
simulations. This would require a self consistent model of the electron cloud
effects like a combination of PyECLOUD and HEADTAIL in one big simulation for
all bunches, which would then allow to study the observed horizontal coupled
bunch instability. Finally it might be interesting to include the effect of
the transverse feedback in HEADTAIL, as this would help to estimate better the
required settings of chromaticity and octupoles to stabilize the beam in a
strong electron-cloud regime.
## 7 ACKNOWLEDGEMENTS
The authors would like to thank D. Valuch for help extracting and analyzing
the turn-by-turn data from the damper pick-ups and feedback setup, G. Arduini
for valuable discussion, and K. Li for support for the HEADTAIL simulations.
## References
* [1] H. Bartosik and W. Höfle, “Analysis of bunch by bunch oscillations with bunch trains at injection into LHC at 25 ns bunch spacing”, CERN-ATS-Note-2012-027 MD (LHC), Geneva 2012.
* [2] G. Iadarola and G. Rumolo, “PyECLOUD and Build Up Simulations at CERN”, these proceedings.
* [3] G. Rumolo et al., “Electron Cloud Effects in the LHC in 2011”, in Proceedings of the LHC Beam Operation Workshop - Evian 2011, Evian 2011.
* [4] G. Rumolo et al., “LHC Experience with Different Bunch Spacings in 2011 (25, 50 and 75ns)”, in Proceedings of Chamonix 2012 workshop on LHC Performance, Chamonix 2012.
* [5] A. Gianluigi et al., “Beam observations with electron cloud in the CERN PS and SPS complex”, 31st Advanced ICFA Beam Dynamics Workshop on Electron-Cloud Effects, Napa 2004.
* [6] G. Rumolo and F. Zimmermann, “Electron cloud simulations: beam instabilities and wakefields”, Phys. Rev. STAccel. Beams 5, 121002, 2002.
* [7] K. Li and G. Rumolo, “Review of Beam Instabilities in the Presence of Electron Clouds in the LHC”, in Proceedings of IPAC11, San Sebastian 2011.
|
arxiv-papers
| 2013-09-25T13:26:37 |
2024-09-04T02:49:51.493017
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "H. Bartosik (CERN), W. Hofle (CERN), G. Iadarola (CERN), Y.\n Papaphilippou (CERN), G. Rumolo (CERN)",
"submitter": "Scientific Information Service CERN",
"url": "https://arxiv.org/abs/1309.6494"
}
|
1309.6516
|
# Prediction of stable hafnium carbides: their stoichiometries, mechanical
properties, and electronic structure
Qingfeng Zeng [email protected] Science and Technology on Thermostructural
Composite Materials Laboratory, School of Materials Science and Engineering,
Northwestern Polytechnical University, Xi’an, Shaanxi 710072, PR China Junhui
Peng Science and Technology on Thermostructural Composite Materials
Laboratory, School of Materials Science and Engineering, Northwestern
Polytechnical University, Xi’an, Shaanxi 710072, PR China Artem R. Oganov
Department of Geosciences, Center for Materials by Design, and Institute for
Advanced Computational Science, State University of New York, Stony Brook, NY
11794-2100, USA Moscow Institute of Physics and Technology, Dolgoprudny,
Moscow Region 141700, Russia School of Materials Science and Engineering,
Northwestern Polytechnical University, Xi’an, Shaanxi 710072, PR China Qiang
Zhu Department of Geosciences, Center for Materials by Design, and Institute
for Advanced Computational Science, State University of New York, Stony Brook,
NY 11794-2100, USA Congwei Xie Science and Technology on Thermostructural
Composite Materials Laboratory, School of Materials Science and Engineering,
Northwestern Polytechnical University, Xi’an, Shaanxi 710072, PR China
Xiaodong Zhang Institute of Modern Physics, Northwest University, Xi’an,
Shaanxi 710069, PR China Science and Technology on Thermostructural Composite
Materials Laboratory, School of Materials Science and Engineering,
Northwestern Polytechnical University, Xi’an, Shaanxi 710072, PR China Dong
Dong Science and Technology on Thermostructural Composite Materials
Laboratory, School of Materials Science and Engineering, Northwestern
Polytechnical University, Xi’an, Shaanxi 710072, PR China Litong Zhang
Science and Technology on Thermostructural Composite Materials Laboratory,
School of Materials Science and Engineering, Northwestern Polytechnical
University, Xi’an, Shaanxi 710072, PR China Laifei Cheng Science and
Technology on Thermostructural Composite Materials Laboratory, School of
Materials Science and Engineering, Northwestern Polytechnical University,
Xi’an, Shaanxi 710072, PR China
###### Abstract
Hafnium carbides are studied by a systematic search for possible stable
stoichiometric compounds in the Hf-C system at ambient pressure using
variable-composition _ab initio_ evolutionary algorithm implemented in the
USPEX code. In addition to well-known HfC, we predicted two additional
compounds Hf3C2 and Hf6C5. The structure of Hf6C5 with space group $C2/m$
contains 11 atoms in the primitive cell and this prediction revives the
earlier proposal by A. I. Gusev. The stable structure of Hf3C2 also has space
group $C2/m$, and is more energetically favorable than the $Immm$,
$P\bar{3}m1$, $P2$ and $C222_{1}$ structures put forward by A. I. Gusev.
Dynamical and mechanical stability of the newly predicted structures have been
verified by calculations of their phonons and elastic constants. The bulk and
shear moduli of Hf3C2 are 195.8 GPa and 143.1 GPa, respectively, while for
Hf6C5 they are 227.9 GPa and 187.2 GPa, respectively. Their mechanical
properties are inferior to those of HfC due to the presence of structural
vacancies. Chemical bonding, band structure, and Bader charge are presented
and discussed.
††preprint: version 0
## I Introduction
Hafnium carbides, known as ultra-high temperature ceramics, have attracted
growing attention because of their unique features. These include extremely
high melting temperature and hardness, high thermal and electrical
conductivity, and chemical stability, and make them promising advanced
materials even in extreme thermal and chemical environments. Levine _et al._
(2002); Opeka _et al._ (2004); Savino _et al._ (2005)
Hafnium carbide is known to crystallize in the NaCl-type structure (space
group $Fm\bar{3}m$) and to have the composition HfC. This is a relatively
well-studied material. Its elastic properties and phonon spectra have been
studied experimentallyBrown _et al._ (1966); Smith and Gläser (1970), and its
structural, elastic, electronic and phonon properties have been computed using
first-principles methodsLi _et al._ (2011a, b). However, according to the
results of theoretical calculation by A. I. Gusev Gusev and Rempel (1993),
ordered stoichiometric phases Hf3C2 and Hf6C5 should exist, with possible
space groups $Immm$, $P\bar{3}m1$, $P2$ or $C222_{1}$ for Hf3C2, and $C2/m$,
$P3_{1}$ or $C2$ for Hf6C5. Experimental synthesis and structure determination
of these subtle ordered states encounters problems due to the lack of direct
methods. Gusev and Zyryanova Gusev and Zyryanova (2000) studied the order-
disorder transition of Hf-C system by measuring the magnetic susceptiblity and
confirmed the existence of Hf3C2 and the possible existence of Hf6C5.
Here we explore the stable compounds in the Hf-C system and their crystal
structures at ambient pressure using the variable-composition _ab initio_
evolutionary algorithm Oganov _et al._ (2010, 2011), and discuss their
structures, elastic properties and chemical bonding. In Section 2, we describe
the computational methods that were used in this work. In Section 3, we
present the results - crystal structures, elastic properties and analysis of
the electronic structure. Section 4 presents conclusions of this study.
## II Computational methodology
The prediction of stable compounds and their crystal structures was performed
using evolutionary algorithm implemented in USPEXOganov and Glass (2006);
Lyakhov _et al._ (2013) code developed by Oganov’s group. This approach
features global optimization with real-space representation and flexible
physically motivated variation operators. For every candidate structure
generated by USPEX, we use first principles structural relaxation, based on
density functional theory (DFT) within the Perdew-Burke-Ernzerhof (PBE)
generalized gradient approximation (GGA)Perdew _et al._ (1996), as
implemented in the VASP code Kresse and Furthmüller (1996). The all-electron
projector-augmented wave (PAW) method Blöchl (1994), with the plane-wave
kinetic energy cutoff of 900 eV and k-point meshes with reciprocal-space
resolution of 2$\pi\times$ 0.07 Å-1 were used. These settings enable excellent
convergences of the energy differences, stress tensors and structural
parameters.
The calculations of phonon dispersion, elastic properties, band structure,
density of state were done using CASTEP codeSegall _et al._ (2002) and
utilizing the modification of the Perdew-Burke-Ernzerhof GGA functional for
solids (PBEsol) Perdew _et al._ (2008). Norm-conserving scheme Hamann _et
al._ (1979) is used to generate pseudopotentials for Hf and C with the choices
of electronic configuration of [Xe]6s25d2, and [He]2s22p2, respectively. In
these plane-wave calculations, the cutoff energy is 750 eV, and k-point
separation is 2$\pi\times$ 0.05 Å-1. Bader charge analysisBader (1990) was
performed using Bader Charge Analysis codeTang _et al._ (2009); Sanville _et
al._ (2007); Henkelman _et al._ (2006).
## III Results and discussions
### III.1 Crystal structure prediction and structural properties
Variable-composition evolutionary algorithmOganov _et al._ (2010, 2011) used
in this work is very effective in simultaneously predicting stable
compositions and their structures for multi-component systems. In our
searches, we allowed all possible compositions in the Hf-C system with
structures containing up to 30 atoms in the unit cell. The initial generation
consisted of 50 structures, and all subsequent generations had 40 structures.
50% of new structures were produced by heredity, 20% by softmutation, 10% by
transmutation, 20% by random symmetric structure generator. Stable structures
and their compositions were determined using the convex hull construction: a
compound is thermodynamically stable if the enthalpy of its decomposition into
any other compounds is positive.
We explored the possible stable crystal structures for Hf-C system at ambient
pressure and zero kelvin. In addition to rocksalt-type (B1) HfC (space group
$Fm\bar{3}m$), we found another two compounds Hf3C2 and Hf6C5, both belonging
to space group $C2/m$. The enthalpies of formation of the predicted structures
are shown in Fig. 1. It can be clearly seen that HfC, Hf6C5 and Hf3C2 are
thermodynamically stable compounds.
Figure 1: Convex hull of Hf-C system at ambient pressure. The circle presents
different structure and those structures locating on lines are
thermodynamically stable.
The crystal structures of Hf3C2 and Hf6C5 are shown in Fig. 2, and their
crystallographic data and enthalpies of formation are listed in Table 1. For
comparison, we also present computational results on HfC in Table 1. Structure
of Hf3C2 has space group of $C2/m$ and 20 atoms in the conventional unit cell
(Fig. 2b). Two more structures of Hf3C2 ($Immm$ and $P\bar{3}m1$) proposed by
A. I. GusevGusev and Rempel (1993) were also found during our searches.
However, their enthalpies are higher than that of $C2/m$, which is therefore
more stable. Crystal structure of Hf6C5 is shown in Fig. 2c. Its space group
is also $C2/m$ and its structure has 22 atoms in the conventional unit cell,
its structural parameters are presented in Table 1, and agree with the
theoretical calculation of A. I. GusevGusev and Rempel (1993).
Figure 2: The crystal structures of (a) HfC with one unit cell, and (b) Hf3C2 (c) Hf6C5 with twice unit cells. The space group of HfC, Hf3C2 and Hf6C5 is $Fm\bar{3}m$,$C2/m$, $C2/m$, respectively. Table 1: Crystallographic data and enthalpies of formation of Hf3C2 and Hf6C5, the data of HfC is also presented as a comparison Compound | Space group | Volume | Lattice constants | Atom position | Enthalpy of formation
---|---|---|---|---|---
| (No.) | (Å3/f.u.) | (Å) | (Wyckoff position) | (eV/atom)
Hf3C2 | $C2/m$ | 75.91 | $a$=5.720 | Hf(4i) (0.737, 0.500, 0.744) | -0.901
| (12) | | $b$=9.893 | Hf(8j) (0.248, 0.161, 0.264) |
| | | $c$=5.701 | C(2a) (0.0, 0.0, 0.0) |
| | | | C(2d) (0.0, 0.5, 0.5) |
| | | | C(4h) (0.0, 0.834, 0.5) |
Hf6C5 | $C2/m$ | 153.14 | $a$=5.729 | Hf(4i) (0.739, 0.000, 0.260) | -0.995
| (12) | | $b$=9.900 | Hf(8j) (0.241, 0.327, 0.746) |
| | | $c$=5.731 | C(4g) (0.0, 0.333, 0.0) |
| | | | C(2d) (0.0, 0.5, 0.5) |
| | | | C(4h) (0.0, 0.832, 0.5) |
HfC | $Fm\bar{3}m$ | 25.34 | $a$=4.675 | Hf(4a) (0.0, 0.00, 0.00) | -1.027
| (225) | | $a$=4.637Li _et al._ (2011b) | C(4b) (0.5, 0.5, 0.5) |
| | | $a$=4.639Nartowski _et al._ (1999) | |
All the stable hafnium carbides are strongly related structures and can be
derived from HfC structure with the highest symmetry. HfC has a structure of
cubic-packing hafnium atoms, and carbon atoms fill all octahedral voids(Fig.
2a), which is an ideal cubic rocksalt-type structure. The octahedra shown in
Fig. 2b and c with green color are empty, i.e. formed by six Hf atoms but
without interstitial C atoms. In Hf3C2 structure, only 2/3 of carbon
octahedral voids are filled (and 1/3 are vacant), and in Hf6C5 5/6 are filled
(and 1/6 are vacant). In both Hf3C2 and Hf6C5, the vacancies appear in every
second octahedral layer, with 1/3 of in-layer octahedra occupied (Hf3C2) or
2/3 of in-layer octahedra occupied (Hf6C5). Ordering of the vacancies in both
cases lowers the symmetry from cubic ($Fm\bar{3}m$) to monoclinic ($C2/m$).
Moreover, due to the vacancies, the coordination number of Hf atoms varies in
different systems: 6 in HfC, 5 in Hf6C5, and 4 in Hf3C2, while in all these
structures carbon atoms invariably had the coordination number 6 (octahedral
coordination). In this way, Hf6C5 and Hf3C2 can be described as defective
rocksalt-type structures.
It is instructive to look at molecular volumes (see Table 1). The volumes per
formula unit (f.u.) of Hf3C2 (75.91 Å3/f.u.), Hf6C5 (153.14 Å3/f.u.) and HfC
(25.536 Å3/f.u.) correspond to practically constant volume per Hf atom
(25.3-25.5 Å3 per Hf atom). Fig. 1 shows that the most prominent stable state
is HfC. And Hf3C2 and especially Hf6C5 will be stable only in a narrow range
of chemical potentials in hafnium-rich conditions. This explains why HfC is
well known from experiment, while Hf3C2 and especially Hf6C5 are more elusive.
To verify the dynamical stability of the newly predicted Hf3C2 and Hf6C5, we
computed their phonon dispersions (Fig. 3). No imaginary phonon frequencies
were found throughout the Brillouin zone, suggesting dynamical stability of
these phases.
Figure 3: Phonon dispersion curves of (a) Hf3C2 (b) Hf6C5 at ambient pressure
### III.2 Elastic properties
The elastic constants of a material describe its response to an applied stress
or, conversely, the stress required to maintain a given deformation, and can
be used to evaluate the elastic properties. These properties are directly
related to mechanical stability. The criteria of mechanical stability of a
monoclinic crystal are as followsCowley (1976):
$K_{2}=\det\left|{C_{ij}}\right|,\;\;i,j\leq
5,\;\;K_{2}>0,\;\;C_{44}C_{66}-C_{46}^{2}>0$ (1)
The calculated elastic constants of Hf3C2 and Hf6C5 at the ground state are
listed in Table 2. It is obvious that these criteria are satisfied, suggesting
that Hf3C2 and Hf6C5 are mechanically stable.
The bulk modulus characterizes the response of a material to volume change,
$B=P/(\Delta V/V)$. $P$ is the applied pressure and $\Delta V$ is the volume
change. The shear modulus characterizes the response to shear deformation,
$G=\tau/\gamma$. $\tau$ is the shear stress and $\gamma$ is shear strain. The
bulk modulus $B$ and shear modulus $G$ can be obtained from elastic constants
Hill (1952), and their values for Hf3C2 and Hf6C5 are presented in Table 2.
For comparison, the values of the elastic properties of HfC are also presented
in Table 2, and one can see a good agreement with the values reported in
literature. The bulk and shear moduli of Hf3C2 and Hf6C5 are lower than those
of HfC. The Pugh’s ratiosPugh (1954) of $G_{H}/B_{H}$ of the three compounds
are larger than 0.57, indicating brittleness of these compounds. According to
Eq. 2Chen _et al._ (2011), the Vickers hardness of Hf3C2, Hf6C5, and HfC is
22.28, 30.91 and 32.95 GPa, respectively - lowering of the hardness from HfC
to Hf6C5 to Hf3C2 is an expected consequence of vacancies.
$H_{V}=2*(k^{2}*G)^{0.585}-3$ (2)
Table 2: Calculated elastic constants $C_{ij}$ , the bulk modulus, shear modulus and hardness (GPa) of Hf3C2, Hf6C5 and HfC at the ground state, and some literature values of HfC Compound | Hf3C2 | Hf6C5 | | HfC | |
---|---|---|---|---|---|---
| | | This work | Calc.He _et al._ (2008) | Expt. Weber (1973) | Expt. Brown _et al._ (1966)
$C_{11}$ | 391 | 448 | 589 | 577 | 500 |
$C_{22}$ | 418 | 471 | | | |
$C_{33}$ | 372 | 470 | | | |
$C_{44}$ | 126 | 181 | 192 | 171 | 180 |
$C_{55}$ | 153 | 213 | | | |
$C_{66}$ | 143 | 199 | | | |
$C_{12}$ | 198 | 118 | 98 | 117 | |
$C_{13}$ | 102 | 117 | | | |
$C_{15}$ | -9.7 | -1.4 | | | |
$C_{23}$ | 91 | 96 | | | |
$C_{25}$ | 12 | 28 | | | |
$C_{35}$ | -0.6 | -21 | | | |
$C_{46}$ | -7.2 | 26 | | | |
$B_{H}$ | 195.8 | 227.9 | 261.7 | 270 | | 242
$G_{H}$ | 143.1 | 187.2 | 212.2 | 230 | | 195
$k$111Pugh’s ratio: $k=G_{H}/B_{H}$ | 0.73 | 0.82 | 0.81 | | |
$H_{V}$ | 22.28 | 30.91 | 32.95 | | |
### III.3 Chemical bonding
The band structures and DOS of these three compounds are shown in Fig. 4. All
the stable hafnium carbides are weak metals, as seen from finite but small DOS
at the Fermi level - 3.28, 2.41, and 0.312 electrons/eV for Hf3C2, Hf6C5 and
HfC, respectively. Taking into account the number of atoms in the primitive
cells of Hf3C2, Hf6C5 and HfC, which were used in the calculations of the band
structure and DOS, we see that the DOS at the Fermi level normalized per
valence electron decreases as the number of vacancies decreases - from Hf3C2
(0.082 states/eV) to Hf6C5 (0.054 states/eV) and to HfC (0.039 states/eV). In
all these three compounds there are pronounced pseudogaps at the Fermi level,
and bonding can be characterized as mixed metallic-covalent. Indeed, orbital-
projected DOS indicates strong hybridization of C-$p$ and Hf-$d$ valence
states below the Fermi energy, i.e. presence of significant covalency in all
three compounds. Bader charge analysis show that each Hf atom gives 1.734
electrons to each C atom in HfC. In case of Hf6C5, Hf atoms contribute
1.521$\pm 0.013$ electrons/atom, and C atoms get 1.826$\pm 0.019$
electrons/atom. In case of Hf3C2, Hf atoms contribute 1.267$\pm 0.026$
electrons/atom, and C atoms get 1.90$\pm 0.044$ electrons/atom. This, in
agreement with the DOS, shows lower metallicity of HfC compared that of Hf6C5
and Hf3C2. Thus, HfC could have the highest hardness and melting point among
these three compounds.
Figure 4: The band structure and density of states of (a) Hf3C2 (b) Hf6C5 (c)
HfC
## IV Conclusions
In this paper we explored the possible stable compounds and structures in the
Hf-C system at ambient pressure using variable-composition evolutionary
algorithm. Besides the well-known HfC ($Fm\bar{3}m$), another two
stoichiometric compounds (Hf3C2, $C2/m$ and Hf6C5, $C2/m$) are found. All
three stable hafnium carbides have rocksalt-type structures: HfC with ideal
cubic structure without vacancies, and defective rocksalt-type phases Hf3C2
and Hf6C5 with monoclinic symmetry ($C2/m$) due to ordering of carbon
vacancies. Their elastic constants and phonon dispersions are also calculated,
which verify their mechanical and dynamical stabilities. Their bulk modulus
and shear modulus are lower than that of HfC due to the presence of structural
vacancies. We found that all three compounds are weak metals, with increasing
metallicity as the concentration of vacancies increases. We also found
significant covalency in all these weakly metallic compounds, while their
hardness values fall with increasing concentration of vacancies.
###### Acknowledgements.
We thank the Research Fund of the State Key Laboratory of Solidification
Processing of NWPU(No. 65-TP-2011), the Basic Research Foundation of NWPU (No.
JCY20130114), the Foreign Talents Introduction and Academic Exchange Program
(No. B08040), the National Science Foundation (Nos. EAR-1114313, DMR-1231586),
DARPA (Nos. W31P4Q1310005, W31P4Q1210008), and the Government of the Russian
Federation (No. 14.A12.31.0003) for financial support. The authors also
acknowledge the High Performance Computing Center of NWPU for the allocation
of computing time on their machines. USPEX code, with options for global
optimization of the thermodynamic potential (energy, enthalpy, free energy),
hardness, bandgap, dielectric constant, and other properties, is available at:
http://uspex.stonybrook.edu.
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|
arxiv-papers
| 2013-09-25T14:21:35 |
2024-09-04T02:49:51.500467
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Qingfeng Zeng, Junhui Peng, Artem R. Oganov, Qiang Zhu, Congwei Xie,\n Xiaodong Zhang, Dong Dong, Litong Zhang and Laifei Cheng",
"submitter": "Qingfeng Zeng",
"url": "https://arxiv.org/abs/1309.6516"
}
|
1309.6530
|
# The role of fast magnetic reconnection in acceleration zones of microquasars
and AGNs
###### Abstract
Fast magnetic reconnection events, like those associated to solar flares, can
be a very powerful mechanism operating at the jet launching region in the core
of microquasars and AGNs. We have recently found that the magnetic power
released by reconnection between the magnetic field lines of the coronal inner
disk region and the lines anchored into the black hole is able to accelerate
relativistic particles through a first-order Fermi process and produce the
observed radio luminosity from both microquasars and low luminous AGNs. We
also found that the observed correlation between the radio luminosity and the
mass of these sources, spanning $10^{9}$ orders of magnitude in mass, is
naturally explained by this process. Furthermore, recent observations of
$\gamma$-ray emission with strong variability from AGNs and galactic compact
sources indicate extremely small emission regions. In this work, assuming that
this high energy emission is probably originated in the same acceleration
zones that produce the radio emission, we have applied the scenario above of
magnetic reconnection-driven acceleration in the magnetized corona around the
accretion disks to investigate the origin of the high energy outcomes from an
extensive number of sources including high and low luminous AGNs,
microquasars, GRBs and the Crab nebula. As in the earlier analysis, we find
that the radio emission of all microquasars and low luminous AGNs are well
explained by our model. On the other hand, we find correlation of our model
with the $\gamma$-ray emission only for microquasars and a few low luminous
AGNs (e.g., M$87$), while none of the high luminous AGNs or GRBs are fitted,
neither in radio (as previously found) nor in $\gamma$ emission. We attribute
the lack of correlation of the $\gamma$ emission for most of the low luminous
AGNs to the fact that this processed emission does not depend only on the
local magnetic field activity around the source accretion disk (as the radio
Synchrotron emission), but also on other environmental factors like the photon
and density fields. We also find that the emission from both the high luminous
AGNs and the GRBs actually anti-correlates with our model indicating that
their emission is processed further out in the jet (as we claimed in [6], and
in consistency with recent findings in [24]). We conclude from these results
that the emission we see from the low luminous AGNs and microquasars comes
from the nuclear region of their sources and therefore, can be driven by
nuclear magnetic activity connected to the source/accretion disk corona, as in
our model. However, in the case of the high luminous AGNs (e.g. BL Lacs) and
GRBs, the nuclear emission is blocked by the surrounding density and photon
fields and, therefore, we can only see the jet emission further out (which has
already lost correlation with the nuclear conditions). This also explains why
high luminous AGNs do not fit the fundamental plane of luminosity source
versus mass correlation.
## 1 Introduction
Galactic and extragalactic objects such as microquasars and AGNs often exhibit
variability and quasi-periodic ejections of matter that may offer important
clues about the physical processes that occur in their inner regions.
A model to explain the origin of these ejections and the associated
Synchrotron flare radio emission, which has characteristics that resemble
those of solar flares, was developed by de Gouveia Dal Pino and Lazarian ([5];
see also [6]) where they invoked a process of fast magnetic reconnection
between the magnetic field lines that arise from the accretion disk and the
lines of the magnetosphere of the central source. In accretion episodes where
the ratio between the effective disk pressure and magnetic pressure decreases
to values smaller than the unity and the accretion rate approaches the
critical Eddington rate, the magnetic reconnection may become very fast and
release large amounts of magnetic energy power. Part of this energy heats the
coronal and the disk gas and part accelerates particles to relativistic
velocities through a first-order Fermi-like process (see also [12], [13], [8])
that results in a Synchrotron radio power-law spectrum compatible with the
observations. Using this model, [5] and [6] also found that the observed
correlation between the radio luminosity and the mass from microquasars to low
luminous AGNs, spanning $10^{9}$ orders of magnitude in mass, is naturally
explained by this process.
In recent years, the very high energy emission (VHE) from AGNs have also
revealed strong variability, with timescales of the order of days (e.g.,
M$87$), which points to extremely compact emission regions (corresponding to
only a few Schwarzschild radii; e.g. [2]). However, the localization of the
emitting zones for several sources is until now unclear. Magnetic reconnection
events occurring close to the black holes (BH) could offer appropriate
conditions for producing particle acceleration and the associated VHE
$\gamma$-ray emission in these sources. Similarly, microquasars are also
expected to emit high-energy $\gamma$-rays owing to their general similarities
to quasars (e.g., [11]). Until now, only one source of this type, namely Cyg
X-$3$, has been unambiguously detected in the GeV $\gamma$-rays, by the Agile
and Fermi observatories ([21], [1]). At TeV energies, only the flux upper
limits are available, in spite of intensive monitoring. There is also some
evidence of sporadic GeV-TeV $\gamma$-ray emission from another source of this
type, Cyg X-$1$.
In the framework of high luminous AGNs and GRBs, Nemen et al. ([24]) have
recently shown that the jets produced by both classes of sources exhibit
nearly the same correlation between the kinetic power carried by accelerated
particles and the $\gamma$-ray luminosity. This is an indication that the bulk
of the high-energy emission comes from the jet region in these sources. On the
other hand, the lack of correlation with low luminous AGNs may be an
indication that the acceleration zones responsible for the high (and also low)
energy emission in these cases are mostly in the core region, rather than
further out in the jet, just like in the model proposed above. Since the
$\gamma$-ray emission is probably originated in the same acceleration zones
that produce the radio emission, in the present work we apply the de Gouveia
Dal Pino et al. scenario above of reconnection-driven acceleration in the
magnetized corona of the accretion disks ([5] and [6]) and investigate if the
high energy outcomes of microquasars and low luminous AGNs can be also
interpreted in the light of this mechanism.
## 2 Rate of magnetic energy released by magnetic reconnection: application
to radio emission
To evaluate the amount of magnetic energy that can be extracted through
violent magnetic reconnection, it is adopted the standard model for the
radiation-dominated accretion disk by [22] and the model by [4] to quantify
the parameters of the corona. Also, it is assumed that the inner radius of the
accretion disk corresponds approximately to the last stable orbit around the
BH ($R_{X}=3R_{S}$, where $R_{S}$ is the Schwarzschild radius). To determine
the accretion rate immediately before an event of violent magnetic
reconnection, it is assumed the equilibrium between the disk gas ram pressure
and the magnetic pressure of the magnetosphere anchored at the BH horizon. It
is assumed further that the intensity of the BH horizon field is of the order
of that of the inner disk. Under these conditions one can show that the
magnetic energy power released during violent magnetic reconnection for
microquasars is approximately given by [6]:
$\dot{W}_{B}\simeq
10^{35}\alpha_{0.5}^{-19/16}\beta_{0.8}^{-9/16}M_{14}^{1/2}l_{100R_{X}}^{11/16}\leavevmode\nobreak\
\leavevmode\nobreak\ erg/s,$ (1)
and for AGNs by [6]:
$\dot{W}_{B}\simeq
10^{40}\alpha_{0.5}^{-19/16}\beta_{0.8}^{-9/16}M_{8}^{1/2}l_{10^{4}R_{X}}^{11/16}\leavevmode\nobreak\
\leavevmode\nobreak\ erg/s,$ (2)
where $\alpha_{0.5}=\alpha/0.5$ is the disk viscosity and
$\beta_{0.8}=\beta/0.8$ is defined as the ratio between the effective disk
pressure and the magnetic pressure. The BH mass $M_{14}=M/14M_{\odot}$ and
$M_{8}=M/10^{8}M_{\odot}$ are parameters suitable for the microquasars and
AGNs, respectively, and the parameters $l_{100R_{X}}=l/100R_{X}$ and
$l_{10^{4}R_{X}}=l/10^{4}R_{X}$ are the scale height of the magnetic
reconnection zone in the corona.
Figure 1 extracted from [6] depicts a synthesis of the magnetic reconnection
scenario for relativistic sources including both microquasars and AGNs. The
diagram shows the calculated magnetic power released in violent magnetic
reconnection events as a function of the central source mass for a suitable
choice of the parameter space. The symbols correspond to the observed radio
luminosities of superluminal components (stars for microquasars, circles and
triangles for the low luminous AGNs, i.e. LINERs and Seyfert galaxies,
respectively, and squares for luminous AGNs).
Figure 1: $\dot{W}_{B}$ versus the BH mass for both microquasars and AGNs. The
stars represent the observed radio luminosities for microquasars. The circles,
triangles and squares are observed radio luminosities of jets at parsec scales
from LINERS, Seyfert galaxies, and luminous AGNs, respectively. The thick bars
correspond to the calculated magnetic reconnection power and encompass the
parameter space that spans $5M_{\odot}\leq M\leq 10^{10}M_{\odot}$,
$0.05\leq\alpha\leq 0.5$, $0.1\leq\beta\leq 1$, and $1R_{S}\leq l\leq
1000R_{S}$ (or $0.3R_{X}\leq l\leq 333R_{X}$), with $1R_{S}\leq l\leq 10R_{S}$
in black, $10R_{S}<l\leq 30R_{S}$ in light gray, and $30R_{S}<l\leq 1000R_{S}$
in dark gray. Extracted from [6].
The diagram indicates that the magnetic power released during violent
reconnection events obeys a correlation that is maintained throughout this
interval, spanning $10^{9}$ orders of magnitude. This correlation implies an
almost linear dependence (in a log-log diagram), which is approximately
independent of the physical properties of the accretion disks of these
sources. Moreover, it is compatible with the so-called “fundamental plane”
obtained empirically, which correlates the radio and X-rays emission of
microquasars and AGNs with the masses of their BH (see [3]). Thus, the model
of [5] provides a simple physical interpretation for the existence of this
empirical correlation as due to coronal magnetic activity in these sources.
The diagram also reveals that the more luminous AGNs do not obey the same
correlation, possibly because the density around the coronal region in these
sources is so high that it “masks” the emission due to the magnetic activity.
The radio emission in these cases is possibly due to regions further out at
the supersonic jet, where it has already expanded enough to become optically
thin and visible and where the relativistic electrons are probably accelerated
in shocks (see [6]). This conclusion in de Gouveia Dal Pino et al. [6] is
compatible with the recent correlation found by Nemem et al. ([24]).
## 3 The same scenario for the $\gamma$-ray emission
Employing the model described above of [6], we selected a preliminary set of
compact sources including high and low luminous AGNs, GRBs and galactic
sources (e.g., the pulsar Crab and the microquasars Cyg-X$1$ and Cyg-X$3$)
with observed $\gamma$ and radio emission and compared their luminosities with
the calculated magnetic reconnection power as a function of their mass.
Figure 2: Magnetic power versus the central source mass. The open blue symbols
correspond to the $\gamma$-ray emission from blazars, low luminous AGNs and
galactic black holes (Cyg-X$1$ and Cyg-X$3$). The blue vertical line and blue
star correspond, respectively, to the $\gamma$ emission from GRBs and the Crab
pulsar. The filled red symbols correspond to the radio emission of different
types of sources (Seyferts, liners, radio galaxies and microquasars).
Figure 2 depicts the preliminary results of this study. As in Figure 1, it
compares the calculated magnetic power released by fast magnetic reconnection
at the corona/accretion disk region as a function of the source mass with the
observed $\gamma$ and radio luminosities. We have considered the $\gamma$
emission (open blue symbols in Figure 2) of an extensive number of high
luminous AGNs and GRBs (all obtained from [24]), low luminous AGNs (Seyferts,
liners and radio galaxies; see [18], [9], [10] and [15]) and microquasars
(Cyg-X$1$ [14] and Cyg-X$3$ [16]). In the case of GRBs (blue vertical line in
the diagram), we have assumed that all sources have a mass around a few tens
of solar masses. We have also included the radio emission counterparts of
several of these sources (filled red symbols; obtained from [3], and [23]) and
the recently detected $\gamma$-ray emission from the Crab pulsar (blue star)
which apparently comes from the surrounds of the central source (see [20]).
## 4 Discussion and Conclusions
As in the previous works ([5] and [6]), we clearly see in Figure 2 that the
observed radio luminosity of the low luminous AGNs and microquasars can be
explained by the magnetic power released by fast reconnection in the core
region of these sources. This emission is due to Synchrotron radiation from
relativistic electrons which can be accelerated by a first-order Fermi process
directly within the magnetic reconnection site in the coronal region around
the accretion disk ([5], [12], [13], [8]). The corresponding $\gamma$-ray
emission from these sources, which is produced from the interaction of the
accelerated relativistic electrons and protons with the surrounding photon and
density fields (through inverse Compton, and/or pp inelastic collisions and
photon-meson decay) can in principle be also associated with the same emission
zone in the surroundings of the core of these sources. However, direct
evidence for this association, though found for the microquasars and the Crab,
is not found for most of the low luminous AGNs. This is explained by the fact
that the $\gamma$ emission (contrary to the radio Synchrotron emission) does
not depend only on the local magnetic fields, but also on the photon and
density fields in the surroundings of the source/accretion disk, as stressed
above, and these factors can provoke the loss of correlation with our nuclear
emission model (induced by magnetic activity around the accretion disk).
The high luminous AGNs and the GRBs in Figure 2, on the other hand, clearly do
not have their radio or $\gamma$ emission correlated with the magnetic
reconnection power released at the core regions. This result confirms the
previous findings of [6] and [24] (see also [7]) which suggested that the
$\gamma$ and radio emission observed in such sources is originated further out
at the relativistic jet associated to these sources (as the nuclear emission
is being screened by the surrounding strong photon and density fields).
The results above connecting both the radio and $\gamma$ emission from low
luminous compact sources to magnetically dominated reconnection processes in
their nuclear regions, though preliminary are very promising as they suggest a
unifying single process of relativistic particle acceleration in the core
region of low luminous AGNs and compact galactic sources which naturally
interpret the Fundamental Plane [3]. In forthcoming work, we intend to extend
the analysis of the diagram above including more radio compact sources with
$\gamma$-ray emission counterparts in order to reinforce the present
conclusions.
Acknowledgment: This work has been partially supported by the Brazilian
funding agencies FAPESP, CNPq and CAPES.
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|
arxiv-papers
| 2013-09-25T14:46:22 |
2024-09-04T02:49:51.508503
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "L.H.S. Kadowaki and E.M. de Gouveia Dal Pino",
"submitter": "Lu\\'is Henrique Sinki Kadowaki",
"url": "https://arxiv.org/abs/1309.6530"
}
|
1309.6534
|
EUROPEAN ORGANIZATION FOR NUCLEAR RESEARCH (CERN)
CERN-PH-EP-2013-176 LHCb-PAPER-2013-053 18 December 2013
Measurement of $D^{0}$–$\kern 4.14793pt\overline{\kern-4.14793ptD}{}^{0}$
mixing parameters and search for $C\\!P$ violation using $D^{0}\rightarrow
K^{+}\pi^{-}$ decays
The LHCb collaboration†††Authors are listed on the following pages.
Measurements of charm mixing parameters from the decay-time-dependent ratio of
$D^{0}\rightarrow K^{+}\pi^{-}$ to $D^{0}\rightarrow K^{-}\pi^{+}$ rates and
the charge-conjugate ratio are reported. The analysis uses data, corresponding
to $3\mbox{\,fb}^{-1}$ of integrated luminosity, from proton-proton collisions
at 7 and 8 $\mathrm{\,Te\kern-1.00006ptV}$ center-of-mass energies recorded by
the LHCb experiment. In the limit of charge-parity ($C\\!P$) symmetry, the
mixing parameters are determined to be $x^{\prime 2}=(5.5\pm 4.9)\times
10^{-5}$, $y^{\prime}=(4.8\pm 1.0)\times 10^{-3}$, and $R_{D}=(3.568\pm
0.066)\times 10^{-3}$. Allowing for $C\\!P$ violation, the mixing parameters
are determined separately for $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ mesons yielding $A_{D}=(-0.7\pm
1.9)\%$, for the direct $C\\!P$-violating asymmetry, and $0.75<|q/p|<1.24$ at
the $68.3\%$ confidence level, where $q$ and $p$ are parameters that describe
the mass eigenstates of the neutral charm mesons in terms of the flavor
eigenstates. This is the most precise determination of these parameters from a
single experiment and shows no evidence for $C\\!P$ violation.
© CERN on behalf of the LHCb collaboration, license CC-BY-3.0.
LHCb collaboration
R. Aaij40, B. Adeva36, M. Adinolfi45, C. Adrover6, A. Affolder51, Z.
Ajaltouni5, J. Albrecht9, F. Alessio37, M. Alexander50, S. Ali40, G.
Alkhazov29, P. Alvarez Cartelle36, A.A. Alves Jr24, S. Amato2, S. Amerio21, Y.
Amhis7, L. Anderlini17,f, J. Anderson39, R. Andreassen56, J.E. Andrews57, R.B.
Appleby53, O. Aquines Gutierrez10, F. Archilli18, A. Artamonov34, M. Artuso58,
E. Aslanides6, G. Auriemma24,m, M. Baalouch5, S. Bachmann11, J.J. Back47, A.
Badalov35, C. Baesso59, V. Balagura30, W. Baldini16, R.J. Barlow53, C.
Barschel37, S. Barsuk7, W. Barter46, Th. Bauer40, A. Bay38, J. Beddow50, F.
Bedeschi22, I. Bediaga1, S. Belogurov30, K. Belous34, I. Belyaev30, E. Ben-
Haim8, G. Bencivenni18, S. Benson49, J. Benton45, A. Berezhnoy31, R. Bernet39,
M.-O. Bettler46, M. van Beuzekom40, A. Bien11, S. Bifani44, T. Bird53, A.
Bizzeti17,h, P.M. Bjørnstad53, T. Blake37, F. Blanc38, J. Blouw10, S. Blusk58,
V. Bocci24, A. Bondar33, N. Bondar29, W. Bonivento15, S. Borghi53, A.
Borgia58, T.J.V. Bowcock51, E. Bowen39, C. Bozzi16, T. Brambach9, J. van den
Brand41, J. Bressieux38, D. Brett53, M. Britsch10, T. Britton58, N.H. Brook45,
H. Brown51, A. Bursche39, G. Busetto21,q, J. Buytaert37, S. Cadeddu15, O.
Callot7, M. Calvi20,j, M. Calvo Gomez35,n, A. Camboni35, P. Campana18,37, D.
Campora Perez37, A. Carbone14,c, G. Carboni23,k, R. Cardinale19,i, A.
Cardini15, H. Carranza-Mejia49, L. Carson52, K. Carvalho Akiba2, G. Casse51,
L. Castillo Garcia37, M. Cattaneo37, Ch. Cauet9, R. Cenci57, M. Charles54, Ph.
Charpentier37, S.-F. Cheung54, N. Chiapolini39, M. Chrzaszcz39,25, K. Ciba37,
X. Cid Vidal37, G. Ciezarek52, P.E.L. Clarke49, M. Clemencic37, H.V. Cliff46,
J. Closier37, C. Coca28, V. Coco40, J. Cogan6, E. Cogneras5, P. Collins37, A.
Comerma-Montells35, A. Contu15,37, A. Cook45, M. Coombes45, S. Coquereau8, G.
Corti37, B. Couturier37, G.A. Cowan49, D.C. Craik47, M. Cruz Torres59, S.
Cunliffe52, R. Currie49, C. D’Ambrosio37, P. David8, P.N.Y. David40, A.
Davis56, I. De Bonis4, K. De Bruyn40, S. De Capua53, M. De Cian11, J.M. De
Miranda1, L. De Paula2, W. De Silva56, P. De Simone18, D. Decamp4, M.
Deckenhoff9, L. Del Buono8, N. Déléage4, D. Derkach54, O. Deschamps5, F.
Dettori41, A. Di Canto11, H. Dijkstra37, M. Dogaru28, S. Donleavy51, F.
Dordei11, P. Dornan52, A. Dosil Suárez36, D. Dossett47, A. Dovbnya42, F.
Dupertuis38, P. Durante37, R. Dzhelyadin34, A. Dziurda25, A. Dzyuba29, S.
Easo48, U. Egede52, V. Egorychev30, S. Eidelman33, D. van Eijk40, S.
Eisenhardt49, U. Eitschberger9, R. Ekelhof9, L. Eklund50,37, I. El Rifai5, Ch.
Elsasser39, A. Falabella14,e, C. Färber11, C. Farinelli40, S. Farry51, D.
Ferguson49, V. Fernandez Albor36, F. Ferreira Rodrigues1, M. Ferro-Luzzi37, S.
Filippov32, M. Fiore16,e, C. Fitzpatrick37, M. Fontana10, F. Fontanelli19,i,
R. Forty37, O. Francisco2, M. Frank37, C. Frei37, M. Frosini17,37,f, E.
Furfaro23,k, A. Gallas Torreira36, D. Galli14,c, M. Gandelman2, P. Gandini58,
Y. Gao3, J. Garofoli58, P. Garosi53, J. Garra Tico46, L. Garrido35, C.
Gaspar37, R. Gauld54, E. Gersabeck11, M. Gersabeck53, T. Gershon47, Ph. Ghez4,
V. Gibson46, L. Giubega28, V.V. Gligorov37, C. Göbel59, D. Golubkov30, A.
Golutvin52,30,37, A. Gomes2, P. Gorbounov30,37, H. Gordon37, M. Grabalosa
Gándara5, R. Graciani Diaz35, L.A. Granado Cardoso37, E. Graugés35, G.
Graziani17, A. Grecu28, E. Greening54, S. Gregson46, P. Griffith44, L.
Grillo11, O. Grünberg60, B. Gui58, E. Gushchin32, Yu. Guz34,37, T. Gys37, C.
Hadjivasiliou58, G. Haefeli38, C. Haen37, S.C. Haines46, S. Hall52, B.
Hamilton57, T. Hampson45, S. Hansmann-Menzemer11, N. Harnew54, S.T. Harnew45,
J. Harrison53, T. Hartmann60, J. He37, T. Head37, V. Heijne40, K. Hennessy51,
P. Henrard5, J.A. Hernando Morata36, E. van Herwijnen37, M. Heß60, A.
Hicheur1, E. Hicks51, D. Hill54, M. Hoballah5, C. Hombach53, W. Hulsbergen40,
P. Hunt54, T. Huse51, N. Hussain54, D. Hutchcroft51, D. Hynds50, V.
Iakovenko43, M. Idzik26, P. Ilten12, R. Jacobsson37, A. Jaeger11, E. Jans40,
P. Jaton38, A. Jawahery57, F. Jing3, M. John54, D. Johnson54, C.R. Jones46, C.
Joram37, B. Jost37, M. Kaballo9, S. Kandybei42, W. Kanso6, M. Karacson37, T.M.
Karbach37, I.R. Kenyon44, T. Ketel41, B. Khanji20, O. Kochebina7, I.
Komarov38, R.F. Koopman41, P. Koppenburg40, M. Korolev31, A. Kozlinskiy40, L.
Kravchuk32, K. Kreplin11, M. Kreps47, G. Krocker11, P. Krokovny33, F. Kruse9,
M. Kucharczyk20,25,37,j, V. Kudryavtsev33, K. Kurek27, T. Kvaratskheliya30,37,
V.N. La Thi38, D. Lacarrere37, G. Lafferty53, A. Lai15, D. Lambert49, R.W.
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Leflat31, J. Lefrançois7, S. Leo22, O. Leroy6, T. Lesiak25, B. Leverington11,
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S. Lohn37, I. Longstaff50, J.H. Lopes2, N. Lopez-March38, H. Lu3, D.
Lucchesi21,q, J. Luisier38, H. Luo49, O. Lupton54, F. Machefert7, I.V.
Machikhiliyan30, F. Maciuc28, O. Maev29,37, S. Malde54, G. Manca15,d, G.
Mancinelli6, J. Maratas5, U. Marconi14, P. Marino22,s, R. Märki38, J. Marks11,
G. Martellotti24, A. Martens8, A. Martín Sánchez7, M. Martinelli40, D.
Martinez Santos41,37, D. Martins Tostes2, A. Martynov31, A. Massafferri1, R.
Matev37, Z. Mathe37, C. Matteuzzi20, E. Maurice6, A. Mazurov16,37,e, J.
McCarthy44, A. McNab53, R. McNulty12, B. McSkelly51, B. Meadows56,54, F.
Meier9, M. Meissner11, M. Merk40, D.A. Milanes8, M.-N. Minard4, J. Molina
Rodriguez59, S. Monteil5, D. Moran53, P. Morawski25, A. Mordà6, M.J.
Morello22,s, R. Mountain58, I. Mous40, F. Muheim49, K. Müller39, R. Muresan28,
B. Muryn26, B. Muster38, P. Naik45, T. Nakada38, R. Nandakumar48, I. Nasteva1,
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Zvyagin37.
1Centro Brasileiro de Pesquisas Físicas (CBPF), Rio de Janeiro, Brazil
2Universidade Federal do Rio de Janeiro (UFRJ), Rio de Janeiro, Brazil
3Center for High Energy Physics, Tsinghua University, Beijing, China
4LAPP, Université de Savoie, CNRS/IN2P3, Annecy-Le-Vieux, France
5Clermont Université, Université Blaise Pascal, CNRS/IN2P3, LPC, Clermont-
Ferrand, France
6CPPM, Aix-Marseille Université, CNRS/IN2P3, Marseille, France
7LAL, Université Paris-Sud, CNRS/IN2P3, Orsay, France
8LPNHE, Université Pierre et Marie Curie, Université Paris Diderot,
CNRS/IN2P3, Paris, France
9Fakultät Physik, Technische Universität Dortmund, Dortmund, Germany
10Max-Planck-Institut für Kernphysik (MPIK), Heidelberg, Germany
11Physikalisches Institut, Ruprecht-Karls-Universität Heidelberg, Heidelberg,
Germany
12School of Physics, University College Dublin, Dublin, Ireland
13Sezione INFN di Bari, Bari, Italy
14Sezione INFN di Bologna, Bologna, Italy
15Sezione INFN di Cagliari, Cagliari, Italy
16Sezione INFN di Ferrara, Ferrara, Italy
17Sezione INFN di Firenze, Firenze, Italy
18Laboratori Nazionali dell’INFN di Frascati, Frascati, Italy
19Sezione INFN di Genova, Genova, Italy
20Sezione INFN di Milano Bicocca, Milano, Italy
21Sezione INFN di Padova, Padova, Italy
22Sezione INFN di Pisa, Pisa, Italy
23Sezione INFN di Roma Tor Vergata, Roma, Italy
24Sezione INFN di Roma La Sapienza, Roma, Italy
25Henryk Niewodniczanski Institute of Nuclear Physics Polish Academy of
Sciences, Kraków, Poland
26AGH - University of Science and Technology, Faculty of Physics and Applied
Computer Science, Kraków, Poland
27National Center for Nuclear Research (NCBJ), Warsaw, Poland
28Horia Hulubei National Institute of Physics and Nuclear Engineering,
Bucharest-Magurele, Romania
29Petersburg Nuclear Physics Institute (PNPI), Gatchina, Russia
30Institute of Theoretical and Experimental Physics (ITEP), Moscow, Russia
31Institute of Nuclear Physics, Moscow State University (SINP MSU), Moscow,
Russia
32Institute for Nuclear Research of the Russian Academy of Sciences (INR RAN),
Moscow, Russia
33Budker Institute of Nuclear Physics (SB RAS) and Novosibirsk State
University, Novosibirsk, Russia
34Institute for High Energy Physics (IHEP), Protvino, Russia
35Universitat de Barcelona, Barcelona, Spain
36Universidad de Santiago de Compostela, Santiago de Compostela, Spain
37European Organization for Nuclear Research (CERN), Geneva, Switzerland
38Ecole Polytechnique Fédérale de Lausanne (EPFL), Lausanne, Switzerland
39Physik-Institut, Universität Zürich, Zürich, Switzerland
40Nikhef National Institute for Subatomic Physics, Amsterdam, The Netherlands
41Nikhef National Institute for Subatomic Physics and VU University Amsterdam,
Amsterdam, The Netherlands
42NSC Kharkiv Institute of Physics and Technology (NSC KIPT), Kharkiv, Ukraine
43Institute for Nuclear Research of the National Academy of Sciences (KINR),
Kyiv, Ukraine
44University of Birmingham, Birmingham, United Kingdom
45H.H. Wills Physics Laboratory, University of Bristol, Bristol, United
Kingdom
46Cavendish Laboratory, University of Cambridge, Cambridge, United Kingdom
47Department of Physics, University of Warwick, Coventry, United Kingdom
48STFC Rutherford Appleton Laboratory, Didcot, United Kingdom
49School of Physics and Astronomy, University of Edinburgh, Edinburgh, United
Kingdom
50School of Physics and Astronomy, University of Glasgow, Glasgow, United
Kingdom
51Oliver Lodge Laboratory, University of Liverpool, Liverpool, United Kingdom
52Imperial College London, London, United Kingdom
53School of Physics and Astronomy, University of Manchester, Manchester,
United Kingdom
54Department of Physics, University of Oxford, Oxford, United Kingdom
55Massachusetts Institute of Technology, Cambridge, MA, United States
56University of Cincinnati, Cincinnati, OH, United States
57University of Maryland, College Park, MD, United States
58Syracuse University, Syracuse, NY, United States
59Pontifícia Universidade Católica do Rio de Janeiro (PUC-Rio), Rio de
Janeiro, Brazil, associated to 2
60Institut für Physik, Universität Rostock, Rostock, Germany, associated to 11
61Celal Bayar University, Manisa, Turkey, associated to 37
aP.N. Lebedev Physical Institute, Russian Academy of Science (LPI RAS),
Moscow, Russia
bUniversità di Bari, Bari, Italy
cUniversità di Bologna, Bologna, Italy
dUniversità di Cagliari, Cagliari, Italy
eUniversità di Ferrara, Ferrara, Italy
fUniversità di Firenze, Firenze, Italy
gUniversità di Urbino, Urbino, Italy
hUniversità di Modena e Reggio Emilia, Modena, Italy
iUniversità di Genova, Genova, Italy
jUniversità di Milano Bicocca, Milano, Italy
kUniversità di Roma Tor Vergata, Roma, Italy
lUniversità di Roma La Sapienza, Roma, Italy
mUniversità della Basilicata, Potenza, Italy
nLIFAELS, La Salle, Universitat Ramon Llull, Barcelona, Spain
oHanoi University of Science, Hanoi, Viet Nam
pInstitute of Physics and Technology, Moscow, Russia
qUniversità di Padova, Padova, Italy
rUniversità di Pisa, Pisa, Italy
sScuola Normale Superiore, Pisa, Italy
Mass eigenstates of neutral charm mesons are linear combinations of flavor
eigenstates $|D_{1,2}\rangle=p|D^{0}\rangle\pm q|\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}\rangle$, where $p$ and $q$ are
complex parameters. This results in $D^{0}$–$\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ oscillation. In the limit of
charge-parity ($C\\!P$) symmetry, the oscillation is characterized by the
difference in mass $\Delta m\equiv m_{2}-m_{1}$ and decay width
$\Delta\Gamma\equiv\Gamma_{2}-\Gamma_{1}$ between the $D$ mass eigenstates.
These differences are usually expressed in terms of the dimensionless mixing
parameters $x\equiv\Delta m/\Gamma$ and $y\equiv\Delta\Gamma/2\Gamma$, where
$\Gamma$ is the average decay width of neutral $D$ mesons. If $C\\!P$ symmetry
is violated, the oscillation rates for mesons produced as $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ can differ, further enriching the
phenomenology. Both short- and long-distance components of the amplitude
contribute to the time evolution of neutral $D$ mesons [1, 2, 3]. Short-
distance amplitudes could include contributions from non-standard-model
particles or interactions, possibly enhancing the average oscillation rate or
the difference between $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ meson rates. The study of $C\\!P$
violation in $D^{0}$ oscillation may lead to an improved understanding of
possible dynamics beyond the standard model [4, 5, 6, 7].
The first evidence for $D^{0}$–$\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ oscillation was reported in 2007
[8, 9]. By 2009, the hypothesis of no oscillation was excluded with
significance in excess of 10 standard deviations [10] by combining results
from different experiments [8, 9, 11, 12, 13, 14, 15, 16, 17]. In 2012, the
LHCb experiment reported the first observation from a single measurement with
greater than 5 standard deviation significance [18], which has been recently
confirmed by the CDF experiment [19].
This Letter reports a search for $C\\!P$ violation in $D^{0}$–$\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ mixing by comparing the decay-
time-dependent ratio of $D^{0}\rightarrow K^{+}\pi^{-}$ to $D^{0}\rightarrow
K^{-}\pi^{+}$ rates with the corresponding ratio for the charge-conjugate
processes. An improved determination of the $C\\!P$-averaged charm mixing
parameters with respect to our previous measurement [18] is also reported. The
analysis uses data corresponding to $1.0\mbox{\,fb}^{-1}$ of integrated
luminosity from $\sqrt{s}=7\,\mathrm{\,Te\kern-1.00006ptV}$ $pp$ collisions
recorded by LHCb during 2011 and $2.0\mbox{\,fb}^{-1}$ from
$\sqrt{s}=8\,\mathrm{\,Te\kern-1.00006ptV}$ collisions recorded during 2012.
The neutral $D$ flavor at production is determined from the charge of the low-
momentum pion $\pi_{\rm s}^{+}$ in the flavor-conserving strong-interaction
decay $D^{*+}\rightarrow D^{0}\pi_{\rm s}^{+}$. The inclusion of charge-
conjugate processes is implicit unless stated otherwise. The
$D^{*+}\rightarrow D^{0}(\rightarrow K^{-}\pi^{+})\pi_{\rm s}^{+}$ process is
denoted as right sign (RS), and $D^{*+}\rightarrow D^{0}(\rightarrow
K^{+}\pi^{-})\pi_{\rm s}^{+}$ is denoted as wrong sign (WS). The RS decay rate
is dominated by a Cabibbo-favored amplitude. The WS rate arises from the
interfering amplitudes of the doubly Cabibbo-suppressed $D^{0}\rightarrow
K^{+}\pi^{-}$ decay and the Cabibbo-favored $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}\rightarrow K^{+}\pi^{-}$ decay
following $D^{0}$–$\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$
oscillation, each of similar magnitude. In the limit of $|x|,|y|\ll 1$, and
assuming negligible $C\\!P$ violation, the time-dependent ratio $R(t)$ of WS-
to-RS decay rates is [1, 2, 3, 4]
$R(t)\approx R_{D}+\sqrt{R_{D}}\ y^{\prime}\ \frac{t}{\tau}+\frac{x^{\prime
2}+y^{\prime 2}}{4}\left(\frac{t}{\tau}\right)^{2},$ (1)
where $t$ is the decay time, $\tau$ is the average $D^{0}$ lifetime, and
$R_{D}$ is the ratio of suppressed-to-favored decay rates. The parameters
$x^{\prime}$ and $y^{\prime}$ depend linearly on the mixing parameters as
$x^{\prime}\equiv x\cos\delta+y\sin\delta$ and $y^{\prime}\equiv
y\cos\delta-x\sin\delta$, where $\delta$ is the strong-phase difference
between the suppressed and favored amplitudes $\mathcal{A}(D^{0}\rightarrow
K^{+}\pi^{-})/\mathcal{A}(\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}\rightarrow
K^{+}\pi^{-})=-\sqrt{R_{D}}e^{-i\delta}$. Allowing for $C\\!P$ violation, the
WS rates $R^{+}(t)$ and $R^{-}(t)$ of initially produced $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ mesons are functions of
independent sets of mixing parameters $(R_{D}^{\pm},\,x^{\prime
2\pm},\,y^{\prime\pm})$. A difference between $R^{+}_{D}$ and $R^{-}_{D}$
arises if the ratio between the magnitudes of suppressed and favored decay
amplitudes is not $C\\!P$ symmetric (direct $C\\!P$ violation). Violation of
$C\\!P$ symmetry either in mixing $|q/p|\neq 1$ or in the interference between
mixing and decay amplitudes $\phi\equiv\arg\left[q\mathcal{A}(\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}\rightarrow
K^{+}\pi^{-})/p\mathcal{A}(D^{0}\rightarrow K^{+}\pi^{-})\right]-\delta\neq 0$
are usually referred to as indirect $C\\!P$ violation and would result in
differences between $(x^{\prime 2+},\,y^{\prime+})$ and $(x^{\prime
2-},\,y^{\prime-})$.
The LHCb detector [20] is a single-arm forward spectrometer covering the
pseudorapidity range $2<\eta<5$, designed for the study of particles
containing $b$ or $c$ quarks. Detector components particularly relevant for
this analysis are the silicon vertex detector, which provides reconstruction
of displaced vertices of $b$\- and $c$-hadron decays; the tracking system,
which measures charged particle momenta with relative uncertainty that varies
from $0.4\%$ at $5\,{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$ to $0.6\%$ at
$100\,{\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$, corresponding to a typical mass
resolution of approximately $8\,{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ for
a two-body charm-meson decay; and the ring-imaging Cherenkov detectors, which
provide kaon-pion discrimination [21]. The magnet polarity is periodically
inverted and approximately equal amounts of data are collected in each
configuration to mitigate the effects of detection asymmetries. The online
event-selection system (trigger) [22] consists of a first-level hardware stage
based on information from the calorimeter and muon systems, followed by a
software high-level trigger.
Events with $D^{*+}$ candidates consistent with being produced at the $pp$
collision point (primary vertex) are selected following Ref. [18]. In
addition, a WS candidate is discarded if resulting from a $D^{0}$ candidate
that, associated with another pion, also forms a RS candidate with
$M(D^{0}\pi_{\rm s}^{+})$ within $3\,{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$
of the known $D^{*+}$ mass. This removes about 15% of the WS background with
negligible signal loss. The two-body $D^{0}\pi_{\rm s}^{+}$ mass
$M(D^{0}\pi_{\rm s}^{+})$ is computed using the known $D^{0}$ and $\pi^{+}$
masses [23] and their reconstructed momenta [18]. In Ref. [18], we used events
selected by the hardware trigger based on hadron calorimeter transverse-energy
depositions that were geometrically matched with signal final-state tracks. In
the present analysis, we distinguish two trigger categories. One category
consists of events that meet the above trigger requirement (triggered-on-
signal, TOS). The other comprises events with candidates failing the track-
calorimeter matching and events selected based on muon hardware triggers
decisions ($\overline{\text{TOS}}$). The two subsamples contribute
approximately equal signal yields with similar purities. However, they require
separate treatment due to their differing kinematic distributions and trigger-
induced biases.
Figure 1: Distribution of $M(D^{0}\pi_{\rm s}^{+})$ for selected (a) right-
sign $D^{0}\rightarrow K^{-}\pi^{+}$ and (b) wrong-sign $D^{0}\rightarrow
K^{+}\pi^{-}$ candidates.
The RS and WS signal yields are determined by fitting the $M(D^{0}\pi_{\rm
s}^{+})$ distribution of $D^{0}$ candidates with reconstructed mass within
$24\,{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ of the known value. The time-
integrated $M(D^{0}\pi_{\rm s}^{+})$ distributions are shown in Fig. 1. The
smooth background is dominated by favored $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}\rightarrow K^{+}\pi^{-}$ decays
associated with random $\pi_{\rm s}^{+}$ candidates. The sample contains
$1.15\times 10^{5}$ ($1.14\times 10^{5}$) signal WS $D^{0}$ ($\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$) decays and approximately 230
times more RS decays. Yield differences between $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ decays are dominated by
differences in charm-anticharm production rates and reconstruction
efficiencies. Each sample is divided into 13 subsamples according to the
candidate’s decay time, and signal yields are determined for each using shape
parametrizations determined from simulation and tuned to data [18]. We assume
that for a given $D^{*}$ meson flavor, the signal shapes are common to WS and
RS decays, while the descriptions of the background can differ. The decay-
time-dependent WS-to-RS yield ratios $R^{+}$ and $R^{-}$ observed in the
$D^{0}$ and $\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ samples,
respectively, and their difference are shown in Fig. 2. These are corrected
for the relative efficiencies for reconstructing $K^{-}\pi^{+}$ and
$K^{+}\pi^{-}$ final states.
Figure 2: Efficiency-corrected ratios of WS-to-RS yields for (a) $D^{*+}$
decays, (b) $D^{*-}$ decays, and (c) their differences as functions of decay
time in units of $D^{0}$ lifetime. Projections of fits allowing for (dashed
line) no $C\\!P$ violation, (dotted line) no direct $C\\!P$ violation, and
(solid line) full $C\\!P$ violation are overlaid. The abscissa of the data
points corresponds to the average decay time over the bin; the error bars
indicate the statistical uncertainties.
The mixing parameters are determined by minimizing a $\chi^{2}$ variable that
includes terms for the difference between the observed and predicted ratios
and for systematic deviations of parameters
$\chi^{2}=\sum_{i}\left[\left(\frac{r^{+}_{i}-\epsilon_{r}^{+}\widetilde{R}_{i}^{+}}{\sigma_{i}^{+}}\right)^{2}+\left(\frac{r^{-}_{i}-\epsilon_{r}^{-}\widetilde{R}_{i}^{-}}{\sigma_{i}^{-}}\right)^{2}\right]+\chi^{2}_{\epsilon}\,+\,\chi^{2}_{B}+\,\chi^{2}_{\rm
p}\,.$ (2)
The measured WS-to-RS yield ratio and its statistical uncertainty in the
decay-time bin $i$ are denoted by $r_{i}^{\pm}$ and $\sigma_{i}^{\pm}$,
respectively. The predicted value for the WS-to-RS yield ratio
$\widetilde{R}_{i}^{\pm}$ corresponds to the time integral over bin $i$ of Eq.
(1) including bin-specific corrections. These account for small biases due to
the decay-time evolution of the approximately $3\%$ fraction of signal
candidates originating from $b$-hadron decays ($\Delta_{B}$) and of the about
$0.5\%$ component of peaking background from RS decays in which both final-
state particles are misidentified ($\Delta_{\rm p}$) [18]. The relative
efficiency $\epsilon_{r}^{\pm}$ accounts for instrumental asymmetries in the
$K\pi$ reconstruction efficiencies, mainly caused by $K^{-}$ mesons having a
larger interaction cross section with matter than $K^{+}$ mesons. These
asymmetries are measured in data to be in the range $0.8$–$1.2\%$ with $0.2\%$
precision and to be independent of decay time. They are derived from the
efficiency ratio
$\epsilon_{r}^{+}=1/\epsilon_{r}^{-}=\epsilon(K^{+}\pi^{-})/\epsilon(K^{-}\pi^{+})$,
obtained from the product of $D^{-}\rightarrow K^{+}\pi^{-}\pi^{-}$ and
$D^{+}\rightarrow K^{0}_{\rm\scriptscriptstyle
S}(\rightarrow\pi^{+}\pi^{-})\pi^{+}$ event yields divided by the product of
the corresponding charge-conjugate decay yields. No $C\\!P$ violation is
expected or experimentally observed [23] in these decays. Asymmetries due to
$C\\!P$ violation in neutral kaons and their interaction cross-sections with
matter are negligible. The $1\%$ asymmetry between $D^{+}$ and $D^{-}$
production rates [24] cancels in this ratio, provided that the kinematic
distributions are consistent across samples. We weight the $D^{-}\rightarrow
K^{+}\pi^{-}\pi^{-}$ events so that their kinematic distributions match those
in the $D^{+}\rightarrow K^{0}_{\rm\scriptscriptstyle S}\pi^{+}$ sample.
Similarly, these samples are weighted as functions of $K\pi$ momentum to match
the RS momentum spectra. The parameters associated with $\Delta_{B}$,
$\Delta_{\rm p}$, and $\epsilon_{r}$ are determined separately for TOS and
$\overline{\text{TOS}}$ subsets and vary independently in the fit within their
Gaussian constraints $\chi^{2}_{B}$, $\chi^{2}_{\rm p}$, and
$\chi^{2}_{\epsilon}$ [18].
To avoid experimenters’ bias in the $C\\!P$ violation parameters, the
measurement technique is finalized by adding arbitrary offsets to the WS-to-RS
yield ratios for the $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ samples, designed to mimic the
effect of different mixing parameters in the two samples. To rule out global
systematic uncertainties not accounted for in Eq. (2), the data are first
integrated over the whole decay-time spectrum and subsequently divided into
statistically independent subsets according to criteria likely to reveal
biases from specific instrumental effects. These include the number of primary
vertices in the events, the $K$ laboratory momentum, the $\pi_{\rm s}$ impact
parameter $\chi^{2}$ with respect to the primary vertex, the $D^{0}$ impact
parameter $\chi^{2}$ with respect to the primary vertex, the magnetic field
orientation, and the hardware trigger category. The variations of the time-
integrated charge asymmetry in WS-to-RS yield ratios are consistent with
statistical fluctuations. Then, we investigate decay-time-dependent biases by
dividing the time-binned sample according to the magnet polarity and the
number of primary vertices per event. In the TOS sample, differences of WS-to-
RS yield ratios as functions of decay time for opposite magnet polarities
yield $\chi^{2}$ values of 12, 17, and 14 (for 12 degrees of freedom), for
events with one, two, and more than two primary vertices, respectively. The
corresponding $\chi^{2}$ values in the $\overline{\text{TOS}}$ sample, 9, 11,
and 8, suggest a systematically better consistency. Hence, the statistical
uncertainty of each of the WS-to-RS ratios in the TOS samples is increased by
a factor of $\sqrt{17/12}$, following Ref. [23]. These scaled uncertainties
are used in all subsequent fits. Independent analyses of the 2011 and 2012
data yield consistent results. The ratio between RS $D^{0}$ to $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ decay rates is independent of
decay time with a $62\%$ $p$ value and a standard deviation of $0.16\%$,
showing no evidence of correlations between particle identification or
reconstruction efficiency and decay time.
Table 1: Results of fits to the data for different hypotheses on the $C\\!P$ symmetry. The reported uncertainties are statistical and systematic, respectively; ndf indicates the number of degrees of freedom. See App. A for fits results including correlation coefficients. Parameter | Value
---|---
Direct and indirect $C\\!P$ violation
$R_{D}^{+}$ [ | $10^{-3}$] | $3.545$ $\pm$ | $0.082$ $\pm$ | $0.048$
$y^{\prime+}$ [ | $10^{-3}$] | $5.1$ $\pm$ | $1.2$ $\pm$ | $0.7$
$x^{\prime 2+}$ [ | $10^{-5}$] | $4.9$ $\pm$ | $6.0$ $\pm$ | $3.6$
$R_{D}^{-}$ [ | $10^{-3}$] | $3.591$ $\pm$ | $0.081$ $\pm$ | $0.048$
$y^{\prime-}$ [ | $10^{-3}$] | $4.5$ $\pm$ | $1.2$ $\pm$ | $0.7$
$x^{\prime 2-}$ [ | $10^{-5}$] | $6.0$ $\pm$ | $5.8$ $\pm$ | $3.6$
$\chi^{2}/\text{ndf}$ | $85.9/98$
No direct $C\\!P$ violation
$R_{D}$ [ | $10^{-3}$] | $3.568$ $\pm$ | $0.058$ $\pm$ | $0.033$
$y^{\prime+}$ [ | $10^{-3}$] | $4.8$ $\pm$ | $0.9$ $\pm$ | $0.6$
$x^{\prime 2+}$ [ | $10^{-5}$] | $6.4$ $\pm$ | $4.7$ $\pm$ | $3.0$
$y^{\prime-}$ [ | $10^{-3}$] | $4.8$ $\pm$ | $0.9$ $\pm$ | $0.6$
$x^{\prime 2-}$ [ | $10^{-5}$] | $4.6$ $\pm$ | $4.6$ $\pm$ | $3.0$
$\chi^{2}/\text{ndf}$ | $86.0/99$
No $C\\!P$ violation
$R_{D}$ [ | $10^{-3}$] | $3.568$ $\pm$ | $0.058$ $\pm$ | $0.033$
$y^{\prime}$ [ | $10^{-3}$] | $4.8$ $\pm$ | $0.8$ $\pm$ | $0.5$
$x^{\prime 2}$ [ | $10^{-5}$] | $5.5$ $\pm$ | $4.2$ $\pm$ | $2.6$
$\chi^{2}/\text{ndf}$ | $86.4/101$
Three fits are performed to the data shown in Fig. 2. The first allows direct
and indirect $C\\!P$ violation; the second allows only indirect $C\\!P$
violation by constraining $R_{D}^{\pm}$ to a common value; and the third is a
$C\\!P$-conserving fit that constrains all mixing parameters to be the same in
the $D^{0}$ and $\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ samples.
The fit results and their projections are shown in Table 1 and Fig. 2,
respectively. Figure 3 shows the central values and confidence regions in the
$(x^{\prime 2},\,y^{\prime})$ plane. For each fit, 104 WS-to-RS ratio data
points are used, corresponding to 13 ranges of decay time, distinguishing
$D^{*+}$ from $D^{*-}$ decays, TOS from $\overline{\text{TOS}}$ decays, and
2011 data from 2012 data. The consistency with the hypothesis of $C\\!P$
symmetry is determined from the change in $\chi^{2}$ between the fit without
and with $C\\!P$ violation, taking into account the difference in number of
degrees of freedom. The resulting $p$ value, for the fit with direct and
indirect (indirect only) $C\\!P$ violation allowed, is $91\%$ ($81\%$),
showing that the data are compatible with $C\\!P$ symmetry.
Figure 3: Two-dimensional confidence regions in the $(x^{\prime
2},y^{\prime})$ plane obtained (a) without any restriction on $C\\!P$
violation, (b) assuming no direct $C\\!P$ violation, and (c) assuming $C\\!P$
conservation. The dashed (solid) curves in (a) and (b) indicate the contours
of the mixing parameters associated with $\kern
1.79997pt\overline{\kern-1.79997ptD}{}^{0}$ ($D^{0}$) decays. The best-fit
value for $\kern 1.79997pt\overline{\kern-1.79997ptD}{}^{0}$ ($D^{0}$) decays
is shown with an open (filled) point. The solid, dashed, and dotted curves in
(c) indicate the contours of $C\\!P$-averaged mixing parameters at 68.3%,
95.5%, and 99.7% confidence level (CL), respectively. The best-fit value is
shown with a point.
The uncertainties incorporate both statistical and systematic contributions,
since all relevant systematic effects depend on the true values of the mixing
parameters, and are thus incorporated into the fit $\chi^{2}$. These include
the uncertainty in the fraction of charm mesons from $b$-hadron decays, and
their bias on the observed decay time; the uncertainty in the fraction of
peaking background; and the uncertainty in the determination of the
instrumental asymmetry. The statistical uncertainty is determined in a
separate fit and used to calculate the systematic component by subtraction in
quadrature.
Direct $C\\!P$ violation would produce a nonzero intercept at $t=0$ in the
efficiency-corrected difference of WS-to-RS yield ratios between $D^{0}$ and
$\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ mesons shown in Fig. 2 (c).
It is parametrized by the asymmetry measured in the first fit
$A_{D}\equiv{(R_{D}^{+}-R_{D}^{-})/(R_{D}^{+}+R_{D}^{-})}=(-0.7\pm 1.9)\%$.
Indirect $C\\!P$ violation results in a time dependence of the efficiency-
corrected difference of yield ratios. The slope observed in Fig. 2 (c) is
about 5% of the individual slopes of Figs. 2 (a) and (b) and is consistent
with zero. From the results of the fit allowing for direct and indirect
$C\\!P$ violation, a likelihood for $|q/p|$ is constructed using the relations
$x^{\prime\pm}=|q/p|^{\pm 1}(x^{\prime}\cos\phi\pm y^{\prime}\sin\phi)$ and
$y^{\prime\pm}=|q/p|^{\pm 1}(y^{\prime}\cos\phi\mp x^{\prime}\sin\phi)$.
Confidence intervals are derived with a likelihood-ratio ordering and assuming
that the correlations are independent of the true values of the mixing
parameters. The magnitude of $q/p$ is determined to be $0.75<|q/p|<1.24$ and
$0.67<|q/p|<1.52$ at the $68.3\%$ and $95.5\%$ confidence levels,
respectively. Significantly more stringent bounds on $|q/p|$ and additional
information on $\phi$ are available by combining the present results with
other measurements [10], in particular when also using theoretical
constraints, such as the relationship $\tan\phi=x(1-|q/p|^{2})/y(1+|q/p|^{2})$
[25, 26], which applies in the limit that direct $C\\!P$ violation is
negligible.
In summary, $D^{0}$–$\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$
oscillation is studied using $D^{*+}\rightarrow D^{0}(\rightarrow
K^{+}\pi^{-})\pi^{+}$ decays reconstructed in the full sample of $pp$
collisions, corresponding to $3\mbox{\,fb}^{-1}$ of integrated luminosity
collected by the LHCb experiment in 2011 and 2012. Assuming $C\\!P$
conservation, the mixing parameters are measured to be $x^{\prime 2}=(5.5\pm
4.9)\times 10^{-5}$, $y^{\prime}=(4.8\pm 1.0)\times 10^{-3}$, and
$R_{D}=(3.568\pm 0.066)\times 10^{-3}$. The observed parameters are consistent
with, $2.5$ times more precise than, and supersede the results based on a
subset of the present data [18]. Studying $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ decays separately shows no
evidence for $C\\!P$ violation and provides the most stringent bounds on the
parameters $A_{D}$ and $|q/p|$ from a single experiment.
## Acknowledgements
We express our gratitude to our colleagues in the CERN accelerator departments
for the excellent performance of the LHC. We thank the technical and
administrative staff at the LHCb institutes. We acknowledge support from CERN
and from the national agencies: CAPES, CNPq, FAPERJ and FINEP (Brazil); NSFC
(China); CNRS/IN2P3 and Region Auvergne (France); BMBF, DFG, HGF and MPG
(Germany); SFI (Ireland); INFN (Italy); FOM and NWO (The Netherlands); SCSR
(Poland); MEN/IFA (Romania); MinES, Rosatom, RFBR and NRC “Kurchatov
Institute” (Russia); MinECo, XuntaGal and GENCAT (Spain); SNSF and SER
(Switzerland); NAS Ukraine (Ukraine); STFC (United Kingdom); NSF (USA). We
also acknowledge the support received from the ERC under FP7. The Tier1
computing centres are supported by IN2P3 (France), KIT and BMBF (Germany),
INFN (Italy), NWO and SURF (The Netherlands), PIC (Spain), GridPP (United
Kingdom). We are thankful for the computing resources put at our disposal by
Yandex LLC (Russia), as well as to the communities behind the multiple open
source software packages that we depend on.
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* [25] Y. Grossman, Y. Nir, and G. Perez, Testing new indirect CP violation, Phys. Rev. Lett. 103 (2009) 071602, arXiv:0904.0305
* [26] A. L. Kagan and M. D. Sokoloff, On indirect CP violation and implications for $D^{0}$–$\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ and $B^{0}_{s}$–$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ mixing, Phys. Rev. D 80 (2009) 076008, arXiv:0907.3917
## Appendix A Supplemental material
Table 2: Detailed fit results. Reported uncertainties and correlation
coefficients include both statistical and systematic sources.
Direct and indirect $C\\!P$ violation
---
Results | Correlations
Parameter | Fit value | $R_{D}^{+}$ | $y^{\prime+}$ | $x^{\prime 2+}$ | $R_{D}^{-}$ | $y^{\prime-}$ | $x^{\prime 2-}$
$R_{D}^{+}$ [ | $10^{-3}$] | $3.545\pm 0.095$ | $1.000$ | $-0.942$ | $\phantom{-}0.862$ | $-0.016$ | $-0.007$ | $\phantom{-}0.006$
$y^{\prime+}$ [ | $10^{-3}$] | $5.1\pm 1.4$ | | $\phantom{-}1.000$ | $-0.968$ | $-0.007$ | $\phantom{-}0.007$ | $-0.007$
$x^{\prime 2+}$ [ | $10^{-5}$] | $4.9\pm 7.0$ | | | $\phantom{-}1.000$ | $\phantom{-}0.005$ | $-0.007$ | $\phantom{-}0.008$
$R_{D}^{-}$ [ | $10^{-3}$] | $3.591\pm 0.094$ | | | | $\phantom{-}1.000$ | $-0.941$ | $\phantom{-}0.858$
$y^{\prime-}$ [ | $10^{-3}$] | $4.5\pm 1.4$ | | | | | $\phantom{-}1.000$ | $-0.966$
$x^{\prime 2-}$ [ | $10^{-5}$] | $6.0\pm 6.8$ | | | | | | $\phantom{-}1.000$
No direct $C\\!P$ violation
---
Results | Correlations
Parameter | Fit value | $R_{D}$ | $y^{\prime+}$ | $x^{\prime 2+}$ | $y^{\prime-}$ | $x^{\prime 2-}$
$R_{D}$ [ | $10^{-3}$] | $3.568\pm 0.066$ | $1.000$ | $-0.894$ | $\phantom{-}0.770$ | $-0.895$ | $\phantom{-}0.772$
$y^{\prime+}$ [ | $10^{-3}$] | $4.8\pm 1.1$ | | $\phantom{-}1.000$ | $-0.949$ | $\phantom{-}0.765$ | $-0.662$
$x^{\prime 2+}$ [ | $10^{-5}$] | $6.4\pm 5.5$ | | | $\phantom{-}1.000$ | $-0.662$ | $\phantom{-}0.574$
$y^{\prime-}$ [ | $10^{-3}$] | $4.8\pm 1.1$ | | | | $\phantom{-}1.000$ | $-0.950$
$x^{\prime 2-}$ [ | $10^{-5}$] | $4.6\pm 5.5$ | | | | | $\phantom{-}1.000$
No $C\\!P$ violation
---
Results | Correlations
Parameter | Fit value | $R_{D}$ | $y^{\prime}$ | $x^{\prime 2}$
$R_{D}$ [ | $10^{-3}$] | $3.568\pm 0.066$ | $1.000$ | $-0.953$ | $\phantom{-}0.869$
$y^{\prime}$ [ | $10^{-3}$] | $4.8\pm 1.0$ | | $\phantom{-}1.000$ | $-0.967$
$x^{\prime 2}$ [ | $10^{-5}$] | $5.5\pm 4.9$ | | | $\phantom{-}1.000$
|
arxiv-papers
| 2013-09-25T14:54:08 |
2024-09-04T02:49:51.515341
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "LHCb collaboration: R. Aaij, B. Adeva, M. Adinolfi, C. Adrover, A.\n Affolder, Z. Ajaltouni, J. Albrecht, F. Alessio, M. Alexander, S. Ali, G.\n Alkhazov, P. Alvarez Cartelle, A.A. Alves Jr, S. Amato, S. Amerio, Y. Amhis,\n L. Anderlini, J. Anderson, R. Andreassen, J.E. Andrews, R.B. Appleby, O.\n Aquines Gutierrez, F. Archilli, A. Artamonov, M. Artuso, E. Aslanides, G.\n Auriemma, M. Baalouch, S. Bachmann, J.J. Back, A. Badalov, C. Baesso, V.\n Balagura, W. Baldini, R.J. Barlow, C. Barschel, S. Barsuk, W. Barter, Th.\n Bauer, A. Bay, J. Beddow, F. Bedeschi, I. Bediaga, S. Belogurov, K. Belous,\n I. Belyaev, E. Ben-Haim, G. Bencivenni, S. Benson, J. Benton, A. Berezhnoy,\n R. Bernet, M.-O. Bettler, M. van Beuzekom, A. Bien, S. Bifani, T. Bird, A.\n Bizzeti, P.M. Bj{\\o}rnstad, T. Blake, F. Blanc, J. Blouw, S. Blusk, V. Bocci,\n A. Bondar, N. Bondar, W. Bonivento, S. Borghi, A. Borgia, T.J.V. Bowcock, E.\n Bowen, C. Bozzi, T. Brambach, J. van den Brand, J. Bressieux, D. Brett, M.\n Britsch, T. Britton, N.H. Brook, H. Brown, A. Bursche, G. Busetto, J.\n Buytaert, S. Cadeddu, O. Callot, M. Calvi, M. Calvo Gomez, A. Camboni, P.\n Campana, D. Campora Perez, A. Carbone, G. Carboni, R. Cardinale, A. Cardini,\n H. Carranza-Mejia, L. Carson, K. Carvalho Akiba, G. Casse, L. Castillo\n Garcia, M. Cattaneo, Ch. Cauet, R. Cenci, M. Charles, Ph. Charpentier, S.-F.\n Cheung, N. Chiapolini, M. Chrzaszcz, K. Ciba, X. Cid Vidal, G. Ciezarek,\n P.E.L. Clarke, M. Clemencic, H.V. Cliff, J. Closier, C. Coca, V. Coco, J.\n Cogan, E. Cogneras, P. Collins, A. Comerma-Montells, A. Contu, A. Cook, M.\n Coombes, S. Coquereau, G. Corti, B. Couturier, G.A. Cowan, D.C. Craik, M.\n Cruz Torres, S. Cunliffe, R. Currie, C. D'Ambrosio, P. David, P.N.Y. David,\n A. Davis, I. De Bonis, K. De Bruyn, S. De Capua, M. De Cian, J.M. De Miranda,\n L. De Paula, W. De Silva, P. De Simone, D. Decamp, M. Deckenhoff, L. Del\n Buono, N. D\\'el\\'eage, D. Derkach, O. Deschamps, F. Dettori, A. Di Canto, H.\n Dijkstra, M. Dogaru, S. Donleavy, F. Dordei, P. Dornan, A. Dosil Su\\'arez, D.\n Dossett, A. Dovbnya, F. Dupertuis, P. Durante, R. Dzhelyadin, A. Dziurda, A.\n Dzyuba, S. Easo, U. Egede, V. Egorychev, S. Eidelman, D. van Eijk, S.\n Eisenhardt, U. Eitschberger, R. Ekelhof, L. Eklund, I. El Rifai, Ch.\n Elsasser, A. Falabella, C. F\\\"arber, C. Farinelli, S. Farry, D. Ferguson, V.\n Fernandez Albor, F. Ferreira Rodrigues, M. Ferro-Luzzi, S. Filippov, M.\n Fiore, C. Fitzpatrick, M. Fontana, F. Fontanelli, R. Forty, O. Francisco, M.\n Frank, C. Frei, M. Frosini, E. Furfaro, A. Gallas Torreira, D. Galli, M.\n Gandelman, P. Gandini, Y. Gao, J. Garofoli, P. Garosi, J. Garra Tico, L.\n Garrido, C. Gaspar, R. Gauld, E. Gersabeck, M. Gersabeck, T. Gershon, Ph.\n Ghez, V. Gibson, L. Giubega, V.V. Gligorov, C. G\\\"obel, D. Golubkov, A.\n Golutvin, A. Gomes, P. Gorbounov, H. Gordon, M. Grabalosa G\\'andara, R.\n Graciani Diaz, L.A. Granado Cardoso, E. Graug\\'es, G. Graziani, A. Grecu, E.\n Greening, S. Gregson, P. Griffith, L. Grillo, O. Gr\\\"unberg, B. Gui, E.\n Gushchin, Yu. Guz, T. Gys, C. Hadjivasiliou, G. Haefeli, C. Haen, S.C.\n Haines, S. Hall, B. Hamilton, T. Hampson, S. Hansmann-Menzemer, N. Harnew,\n S.T. Harnew, J. Harrison, T. Hartmann, J. He, T. Head, V. Heijne, K.\n Hennessy, P. Henrard, J.A. Hernando Morata, E. van Herwijnen, M. He\\ss, A.\n Hicheur, E. Hicks, D. Hill, M. Hoballah, C. Hombach, W. Hulsbergen, P. Hunt,\n T. Huse, N. Hussain, D. Hutchcroft, D. Hynds, V. Iakovenko, M. Idzik, P.\n Ilten, R. Jacobsson, A. Jaeger, E. Jans, P. Jaton, A. Jawahery, F. Jing, M.\n John, D. Johnson, C.R. Jones, C. Joram, B. Jost, M. Kaballo, S. Kandybei, W.\n Kanso, M. Karacson, T.M. Karbach, I.R. Kenyon, T. Ketel, B. Khanji, O.\n Kochebina, I. Komarov, R.F. Koopman, P. Koppenburg, M. Korolev, A.\n Kozlinskiy, L. Kravchuk, K. Kreplin, M. Kreps, G. Krocker, P. Krokovny, F.\n Kruse, M. Kucharczyk, V. Kudryavtsev, K. Kurek, T. Kvaratskheliya, V.N. La\n Thi, D. Lacarrere, G. Lafferty, A. Lai, D. Lambert, R.W. Lambert, E.\n Lanciotti, G. Lanfranchi, C. Langenbruch, T. Latham, C. Lazzeroni, R. Le Gac,\n J. van Leerdam, J.-P. Lees, R. Lef\\`evre, A. Leflat, J. Lefran\\c{c}ois, S.\n Leo, O. Leroy, T. Lesiak, B. Leverington, Y. Li, L. Li Gioi, M. Liles, R.\n Lindner, C. Linn, B. Liu, G. Liu, S. Lohn, I. Longstaff, J.H. Lopes, N.\n Lopez-March, H. Lu, D. Lucchesi, J. Luisier, H. Luo, O. Lupton, F. Machefert,\n I.V. Machikhiliyan, F. Maciuc, O. Maev, S. Malde, G. Manca, G. Mancinelli, J.\n Maratas, U. Marconi, P. Marino, R. M\\\"arki, J. Marks, G. Martellotti, A.\n Martens, A. Mart\\'in S\\'anchez, M. Martinelli, D. Martinez Santos, D. Martins\n Tostes, A. Martynov, A. Massafferri, R. Matev, Z. Mathe, C. Matteuzzi, E.\n Maurice, A. Mazurov, J. McCarthy, A. McNab, R. McNulty, B. McSkelly, B.\n Meadows, F. Meier, M. Meissner, M. Merk, D.A. Milanes, M.-N. Minard, J.\n Molina Rodriguez, S. Monteil, D. Moran, P. Morawski, A. Mord\\`a, M.J.\n Morello, R. Mountain, I. Mous, F. Muheim, K. M\\\"uller, R. Muresan, B. Muryn,\n B. Muster, P. Naik, T. Nakada, R. Nandakumar, I. Nasteva, M. Needham, S.\n Neubert, N. Neufeld, A.D. Nguyen, T.D. Nguyen, C. Nguyen-Mau, M. Nicol, V.\n Niess, R. Niet, N. Nikitin, T. Nikodem, A. Nomerotski, A. Novoselov, A.\n Oblakowska-Mucha, V. Obraztsov, S. Oggero, S. Ogilvy, O. Okhrimenko, R.\n Oldeman, M. Orlandea, J.M. Otalora Goicochea, P. Owen, A. Oyanguren, B.K.\n Pal, A. Palano, M. Palutan, J. Panman, A. Papanestis, M. Pappagallo, C.\n Parkes, C.J. Parkinson, G. Passaleva, G.D. Patel, M. Patel, G.N. Patrick, C.\n Patrignani, C. Pavel-Nicorescu, A. Pazos Alvarez, A. Pearce, A. Pellegrino,\n G. Penso, M. Pepe Altarelli, S. Perazzini, E. Perez Trigo, A. P\\'erez-Calero\n Yzquierdo, P. Perret, M. Perrin-Terrin, L. Pescatore, E. Pesen, G. Pessina,\n K. Petridis, A. Petrolini, A. Phan, E. Picatoste Olloqui, B. Pietrzyk, T.\n Pila\\v{r}, D. Pinci, S. Playfer, M. Plo Casasus, F. Polci, G. Polok, A.\n Poluektov, E. Polycarpo, A. Popov, D. Popov, B. Popovici, C. Potterat, A.\n Powell, J. Prisciandaro, A. Pritchard, C. Prouve, V. Pugatch, A. Puig\n Navarro, G. Punzi, W. Qian, B. Rachwal, J.H. Rademacker, B.\n Rakotomiaramanana, M.S. Rangel, I. Raniuk, N. Rauschmayr, G. Raven, S.\n Redford, S. Reichert, M.M. Reid, A.C. dos Reis, S. Ricciardi, A. Richards, K.\n Rinnert, V. Rives Molina, D.A. Roa Romero, P. Robbe, D.A. Roberts, A.B.\n Rodrigues, E. Rodrigues, P. Rodriguez Perez, S. Roiser, V. Romanovsky, A.\n Romero Vidal, M. Rotondo, J. Rouvinet, T. Ruf, F. Ruffini, H. Ruiz, P. Ruiz\n Valls, G. Sabatino, J.J. Saborido Silva, N. Sagidova, P. Sail, B. Saitta, V.\n Salustino Guimaraes, B. Sanmartin Sedes, R. Santacesaria, C. Santamarina\n Rios, E. Santovetti, M. Sapunov, A. Sarti, C. Satriano, A. Satta, M. Savrie,\n D. Savrina, M. Schiller, H. Schindler, M. Schlupp, M. Schmelling, B. Schmidt,\n O. Schneider, A. Schopper, M.-H. Schune, R. Schwemmer, B. Sciascia, A.\n Sciubba, M. Seco, A. Semennikov, K. Senderowska, I. Sepp, N. Serra, J.\n Serrano, P. Seyfert, M. Shapkin, I. Shapoval, Y. Shcheglov, T. Shears, L.\n Shekhtman, O. Shevchenko, V. Shevchenko, A. Shires, R. Silva Coutinho, M.\n Sirendi, N. Skidmore, T. Skwarnicki, N.A. Smith, E. Smith, E. Smith, J.\n Smith, M. Smith, M.D. Sokoloff, F.J.P. Soler, F. Soomro, D. Souza, B. Souza\n De Paula, B. Spaan, A. Sparkes, P. Spradlin, F. Stagni, S. Stahl, O.\n Steinkamp, S. Stevenson, S. Stoica, S. Stone, B. Storaci, M. Straticiuc, U.\n Straumann, V.K. Subbiah, L. Sun, W. Sutcliffe, S. Swientek, V. Syropoulos, M.\n Szczekowski, P. Szczypka, D. Szilard, T. Szumlak, S. T'Jampens, M. Teklishyn,\n E. Teodorescu, F. Teubert, C. Thomas, E. Thomas, J. van Tilburg, V.\n Tisserand, M. Tobin, S. Tolk, D. Tonelli, S. Topp-Joergensen, N. Torr, E.\n Tournefier, S. Tourneur, M.T. Tran, M. Tresch, A. Tsaregorodtsev, P.\n Tsopelas, N. Tuning, M. Ubeda Garcia, A. Ukleja, A. Ustyuzhanin, U. Uwer, V.\n Vagnoni, G. Valenti, A. Vallier, R. Vazquez Gomez, P. Vazquez Regueiro, C.\n V\\'azquez Sierra, S. Vecchi, J.J. Velthuis, M. Veltri, G. Veneziano, M.\n Vesterinen, B. Viaud, D. Vieira, X. Vilasis-Cardona, A. Vollhardt, D.\n Volyanskyy, D. Voong, A. Vorobyev, V. Vorobyev, C. Vo\\ss, H. Voss, R. Waldi,\n C. Wallace, R. Wallace, S. Wandernoth, J. Wang, D.R. Ward, N.K. Watson, A.D.\n Webber, D. Websdale, M. Whitehead, J. Wicht, J. Wiechczynski, D. Wiedner, L.\n Wiggers, G. Wilkinson, M.P. Williams, M. Williams, F.F. Wilson, J. Wimberley,\n J. Wishahi, W. Wislicki, M. Witek, G. Wormser, S.A. Wotton, S. Wright, S. Wu,\n K. Wyllie, Y. Xie, Z. Xing, Z. Yang, X. Yuan, O. Yushchenko, M. Zangoli, M.\n Zavertyaev, F. Zhang, L. Zhang, W.C. Zhang, Y. Zhang, A. Zhelezov, A.\n Zhokhov, L. Zhong, A. Zvyagin",
"submitter": "Angelo Di Canto",
"url": "https://arxiv.org/abs/1309.6534"
}
|
1309.6700
|
# Spectral conditions for the existence of specified paths and cycles in
graphs ††thanks: Supported by the NNSF of China (No.11101057), China
Postdoctoral Science Foundation (No.20110491443) and Zhejiang Provincial
Natural Science Foundation of China(LY12A01016).
Mingqing Zhaia , Huiqiu Linb , Shicai Gongc
aSchool of Mathematical Science, Chuzhou University, Anhui, Chuzhou, 239012,
China
bDepartment of Mathematics, East China University of Science and Technology,
Shanghai, 200237, China
cSchool of Science, Zhejiang A & F University, Zhejiang, Linan, 311300, China
Corresponding author. E-mail addresses: [email protected] (M.Zhai);
[email protected] (H.Lin); [email protected](S.Gong).
###### Abstract
Let $G$ be a graph with $n$ vertices and $\lambda_{n}(G)$ be the least
eigenvalue of its adjacency matrix of $G$. In this paper, we give sharp bounds
on the least eigenvalue of graphs without given pathes or cycles and determine
the extremal graphs. This result gives spectral conditions for the existence
of specified paths and cycles in graphs.
AMS Classification: 05C50; 05C38
Keywords: Path; Cycle; Spectral radius; Least eigenvalue; Adjacency matrix
## 1 Introduction
Let $G=(V,E)$ be the graph with vertex set $V$ and edge set $E$. Let $|G|$ and
$e(G)=|E(G)|$ be the order and the size of $G$, respectively. A graph $G$ is
said $H$-free, if $G$ does not contain a copy of $H$. Denote by $P_{t}$ and
$C_{t}$ the path and cycle of order $t$, respectively. Let $N^{d}(u)=\\{v|v\in
V(G),d_{G}(v,u)=d\\}$, where $V(G)$ is the vertex set of $G$ and $d_{G}(v,u)$
is the distance between $u$ and $v$. Particularly, denote $N(u)=N^{1}(u)$ and
$d(u)=|N(u)|$. For two disjoint vertex sets $V_{1}$ and $V_{2}$, let
$e(V_{1},V_{2})$ be the number of edges with one endpoint in $V_{1}$ and
another in $V_{2}$. Let $A(G)$ be the adjacency matrix of a graph $G$. The
largest modulus of an eigenvalue of $A(G)$ is called the spectral radius of
$G$ and denoted by $\rho(G)$. It is known that for a connected graph $G$,
there exists a positive unit eigenvector corresponding to $\rho(G)$, which is
called the Perron vector of $G$.
A central problem of extremal graph theory is as follows: for a given graph
$H$, what is the maximum number of edges of an $H$-free graph of order $n$? In
the past decades, much attention has been paid to the spectral version of
above question, that is, what is the maximum spectral radius of an $H$-free
graph of order $n$? It is known that if $H$ is a star $K_{1,r+1}$, then the
extremal graph is an $r$-regular graph (see [8]). In recent years, Fiedler and
Nikiforov [6] solved the question for $H\cong C_{n}$ or $H\cong P_{n}$.
Nikiforov solved the case $H\cong C_{2k+1}$ for large enough $n$ (see [9]) and
the case that $H$ is a complete graph $K_{r+1}$ or a complete bipartite graph
$K_{2,r+1}$ (see [10]). In [11], Nikiforov showed if $G$ is a $C_{4}$-free
graph with $n$ vertices, then $\rho^{2}(G)-\rho(G)-(n-1)\leq 0$ with equality
if and only if $G$ is the friendship graph (now $n$ is odd). Moreover, he
conjectured if $G$ is a $C_{4}$-free graph with even vertices, then
$\rho^{3}(G)-\rho^{2}(G)-(n-1)\rho(G)+1\leq 0$ with equality if and only if
$G$ is a star of order $n$ with $\frac{n}{2}-1$ disjoint additional edges.
Zhai and Wang [18] showed this conjecture is true. In [12], Nikiforov obtained
the extremal graph for $H\cong P_{2k+2}$ or $H\cong P_{2k+3}$ and gave a
conjecture on $H\cong C_{2k+2}$ (see [12], Conjecture 15.). This paper solves
these two questions in the class of bipartite graphs (see Theorems 2.1 and
2.2).
Let $\lambda_{n}(G)$ be the least eigenvalue of a graph $G$ of order $n$. It
is known that $\lambda_{n}(G)=-\rho(G)$ for a bipartite graph $G$ (see [4]).
Recently, researchers have begun to pay attention to the least eigenvalues of
graphs with a given value of some well-known integer graph invariant: for
instance: order and size [1, 2, 5, 13], unicyclic graphs with a given number
of pendant vertices [7], matching number and independence number [14], number
of cut vertices [15], connectivity, chromatic number [16], domination number
[17]. This paper also gives a spectral extremal characterization on the least
eigenvalue of graphs (see Theorems 1.3 and 1.4).
###### Lemma 1.1
Let $G=<X,Y>$ be a bipartite graph, where $|X|\geq k$, $|Y|\geq k-1$. If $G$
does not contain a copy of $P_{2k+1}$ with both endpoints in $X$, then
$e(G)\leq(k-1)|X|+k|Y|-k(k-1).$ $None$
Equality holds if and only if $G$ is isomorphic to a complete bipartite graph
$K_{|X|,|Y|}$, where $|X|=k$ or $|Y|=k-1$.
###### Lemma 1.2
Let $G$ be a $P_{2k+3}$-free bipartite graph. If $G$ does not contain a
connected component isomorphic to $K_{k+1,k+1}$, then for an arbitrary vertex
$u\in V(G)$,
$e(N(u),N^{2}(u))\leq(k-1)|N(u)|+k|N^{2}(u)|-k(k-1),$ $None$
and equality holds if and only if the subgraph induced by $N(u)\cup N^{2}(u)$
is isomorphic to one of the following:
(i) $K_{|N(u)|,|N^{2}(u)|}$, where $|N(u)|=k$ or $|N^{2}(u)|=k-1$;
(ii) a graph obtained from $K_{k,|N^{2}(u)|}$ and an additional vertex $v\in
N(u)$ by joining it with $k-1$ vertices of $N^{2}(u)$.
###### Theorem 1.3
Let $G$ be a $C_{t}$-free graph of order $n$.
(i) If $t$ is odd or $n<t$, then
$\lambda_{n}(G)\geq-\sqrt{\lfloor\frac{n}{2}\rfloor\lceil\frac{n}{2}\rceil}$,
and equality holds if and only if $G\cong
K_{\lfloor\frac{n}{2}\rfloor,\lceil\frac{n}{2}\rceil}$.
(ii) If $t=2k+2$ and $n\geq t$, then $\lambda_{n}(G)\geq-\sqrt{k(n-k)}$, and
equality holds if and only if $G\cong K_{k,n-k}$.
###### Theorem 1.4
Let $G$ be a $P_{t}$-free graph of order $n$.
(i) If $n<t$, then
$\lambda_{n}(G)\geq-\sqrt{\lfloor\frac{n}{2}\rfloor\lceil\frac{n}{2}\rceil}$.
Equality holds if and only if $G\cong
K_{\lfloor\frac{n}{2}\rfloor,\lceil\frac{n}{2}\rceil}$.
(ii) If $n\geq t$, then $\lambda_{n}(G)\geq-\sqrt{k(n-k)}$, where
$k=\lfloor\frac{t-2}{2}\rfloor$. Equality holds if and only if $G\cong
K_{k,n-k}$ or $G\cong K_{2,2}\cup K_{1}$ for $(n,t)=(5,5)$.
In the remaining part of this paper, we give the proofs of Lemmas 1.1 and 1.2,
and Theorems 1.3 and 1.4.
## 2 Proofs
Proof of Lemma 1.1. We use induction on $|Y|$. If $|Y|=k-1$, then
$e(G)\leq e(K_{|X|,k-1})=(k-1)|X|.$
So, (1) follows and the equality holds if and only if $G\cong K_{|X|,k-1}$.
Now we suppose that the inequality holds for $|Y|\leq M$, where $M\geq k$.
Case 1. There exists a vertex $u\in Y$ such that $d_{G}(u)\leq k-1$, then by
the hypothesis,
$e(G-u)\leq(k-1)|X|+k(|Y|-1)-k(k-1).$
Hence we have
$e(G)=e(G-u)+d_{G}(u)<(k-1)|X|+k|Y|-k(k-1).$
Case 2. For each vertex $u\in Y$, $d_{G}(u)\geq k$. We use induction on $|X|$.
If $|X|=k$, then
$e(G)\leq e(K_{k,|Y|})=k|Y|,$
so (1) follows and the equality holds if and only if $G\cong K_{k,|Y|}$.
Suppose that the inequality holds for $|X|\leq M$ where $M\geq k+1$.
Subcase 2.1. There exists a vertex $v\in X$, such that $d_{G}(v)\leq k-2$,
then by the hypothesis, we have
$e(G-v)\leq(k-1)(|X|-1)+k|Y|-k(k-1).$
Hence
$e(G)=e(G-u)+d_{G}(u)<(k-1)|X|+k|Y|-k(k-1).$
Subcase 2.2. For each vertex $v\in X$, $d_{G}(v)\geq k-1$. We may assume that
$G$ is connected. Otherwise, we use the induction hypothesis on each component
of $G$ and the inequality follows immediately. Now we suppose that $P$ is a
longest path in $G$ with two endpoints $w_{1},w_{2}$.
Subcase 2.2.1. Either $w_{1}$ or $w_{2}$ in $Y$, say $w_{1}\in Y$. Then
$|P\cap Y|\geq|P\cap X|$. Note that $G$ is $P_{2k+2}$-free, then $|P\cap
X|\leq k$. Since $d_{G}(w_{1})\geq k$, thus we have $|P\cap X|=k$ and
$N_{G}(w_{1})=P\cap X$. One can find now each vertex in $P\cap Y$ is an
endpoint of a path of the same length with $P$. Similarly, we have
$N_{G}(w)=P\cap X$ for any $w\in P\cap Y$. Then $X=P\cap X$ (otherwise we can
find a path with length longer than $|P|$ which is forbidden). Thus we have
$d_{G}(u)=k$ for each $u\in Y$ since $d_{G}(u)\geq k$ and $|X|=|P\cap X|=k$.
Therefore $G\cong K_{k,|Y|}$ and then the inequality in (1) is an equality.
Subcase 2.2.2. Both $w_{1}$ and $w_{2}$ are in $X$. Then $|P\cap X|=|P\cap
Y|+1$. Note that $G$ does not contain a copy of $P_{2k+1}$ with both endpoints
in $X$, then $|P\cap Y|\leq k-1$. Since $P$ is a longest path in $G$, thus
$N_{G}(w_{1})\subset P\cap Y$. Hence we have $N_{G}(w_{1})=P\cap Y$ and
$d_{G}(w_{1})=|P\cap Y|=k-1$ since $d_{G}(w_{1})\geq k-1$. Further, each
vertex $u\in X$ is an endpoint of a path of the same length with $P$. So we
have $d_{G}(u)=k-1$ for each vertex $u\in X$. Similarly, we have $Y=P\cap Y$.
Therefore, $G\cong K_{|X|,k-1}$ and then the inequality in (1) is an equality.
Thus we complete the proof. $\Box$
Proof of Lemma 1.2. Denote by $H$ the subgraph of $G$ induced by $N(u)\cup
N^{2}(u)$. By Lemma 1.1, if $H$ does not contain a copy of $P_{2k+1}$ with
both endpoints in $N(u)$, then Inequality (2) holds, and equality holds if and
only if $G\cong K_{|N(u)|,|N^{2}(u)|}$ with $|N(u)|=k$ or $|N^{2}(u)|=k-1$.
Now suppose that $H$ contains a path $P$ of order $2k+1$ with both endpoints
in $N(u)$.
Since $G$ is $P_{2k+3}$-free, $N(u)\backslash V(P)=\emptyset$ which implies
that $|N(u)|=k+1$ and $|N^{2}(u)|\geq k$. We distinguish the following two
cases.
Case 1. There is a vertex $v\in N(u)$ such that $d_{H}(v)\leq k-1$. Now,
$e(H-v)\leq
e(K_{k,|N^{2}(u)|})=k|N^{2}(u)|=(k-1)|N(u)|+k|N^{2}(u)|-(k+1)(k-1),$
and equality holds if and only if $H-v\cong K_{k,|N^{2}(u)|}$. Therefore,
$e(H)=e(H-v)+d_{H}(v)\leq(k-1)|N(u)|+k|N^{2}(u)|-k(k-1)$
If equality holds, then $d_{H}(v)=k-1$.
Case 2. For each vertex $v\in N(u)$, $d_{H}(v)\geq k$. Now let $w\in N(u)$ be
an endpoint of $P$. Note that $P$ is a longest path in $H$. Then $d_{H}(w)\geq
k$ and $N_{H}(w)\subseteq V(P)\cap N^{2}(u)$. However, $|V(P)\cap
N^{2}(u)|=k$. Hence $N_{H}(w)=V(P)\cap N^{2}(u)$. Further, we can find that
each vertex $x\in N(u)$ must be an endpoint of a path of order $2k+1$ in $H$.
Correspondingly, $N_{H}(x)=V(P)\cap N^{2}(u)$. This implies that $V(P)$
induces a copy of $K_{k+1,k}$. Since $G$ is $P_{2k+3}$-free,
$N^{2}(u)\backslash V(P)=\emptyset$. Thus $H\cong K_{k+1,k}$. Since $G$ does
not contain a connected component isomorphic to $K_{k+1,k+1}$,
$N^{3}(u)\neq\emptyset$. We then have a path of order $2k+3$ in $G$, a
contradiction.
This completes the proof. $\Box$
###### Theorem 2.1
Let $G$ be a $C_{t}$-free bipartite graph of order $n$.
(i) If $t$ is odd or $n<t$, then
$\rho(G)\leq\sqrt{\lfloor\frac{n}{2}\rfloor\lceil\frac{n}{2}\rceil}$. Equality
holds if and only if $G\cong
K_{\lfloor\frac{n}{2}\rfloor,\lceil\frac{n}{2}\rceil}$.
(ii) If $t=2k+2$ and $n\geq t$, then $\rho(G)\leq\sqrt{k(n-k)}$. Equality
holds if and only if $G\cong K_{k,n-k}$.
Proof. (i) Recall that $K_{\lfloor\frac{n}{2}\rfloor,\lceil\frac{n}{2}\rceil}$
has maximal spectral radius among all bipartite graphs of order $n$. If $t$ is
odd or $n<t$, then
$\rho(G)\leq\sqrt{\lfloor\frac{n}{2}\rfloor\lceil\frac{n}{2}\rceil}$ since $G$
is bipartite. Equality holds if and only if $G\cong
K_{\lfloor\frac{n}{2}\rfloor,\lceil\frac{n}{2}\rceil}$.
(ii) Let $G$ be a bipartite graph without $C_{2k+2}$ and
$B=(b_{ij})_{n\times n}=A(G)^{2}-k(n-k)I,$
where $I$ is the identity matrix of order $n$. Since $G$ is $C_{2k+2}$-free,
$G$ does not contain a copy of $P_{2k+1}$ with both endpoints in $N(u)$ for
any $u\in V(G)$. Let $b(u)$ be the sum of the $u$-th row of $B$. Note that
$(A(G)^{2})_{u}$ is the number of walks of length 2 which are started from
$u$. Then by Lemma 1.1, for any $u\in V(G)$,
$\displaystyle b(u)$ $\displaystyle=$
$\displaystyle|N(u)|+e(N(u),N^{2}(u))-k(n-k)$ $\displaystyle\leq$
$\displaystyle|N(u)|+(k-1)|N(u)|+k|N^{2}(u)|-k(k-1)-k(n-k)$ $\displaystyle=$
$\displaystyle k(|N(u)|+|N^{2}(u)|)-k(n-1)$ $\displaystyle\leq$ $\displaystyle
k(n-1)-k(n-1)$ $\displaystyle=$ $\displaystyle 0.$
If the equality holds, then by the second inequality, $N^{3}(u)=\emptyset$.
Further, by the first inequality and Lemma 1.1, $G-u\cong
K_{|N(u)|,|N^{2}(u)|}$ with $|N(u)|=k$ or $|N^{2}(u)|=k-1$. This implies that
$G\cong K_{k,n-k}.$
Let $X$ be an eigenvector of $G$ corresponding to $\rho(G)$ with
$\sum_{i=1}^{n}x_{i}=1$. Then
$\rho^{2}(G)-k(n-k)=\sum_{i=1}^{n}[\rho^{2}(G)-k(n-k)]x_{i}=\sum\limits_{i=1}^{n}\sum\limits_{j=1}^{n}b_{ij}x_{j}=\sum\limits_{u\in
V(G)}b(u)x_{u}\leq 0.$
Thus we have $\rho(G)\leq\sqrt{k(n-k)}$. If the equality holds, then $b(u)=0$
for any $u\in V(G)$. Thus $G\cong K_{k,n-k}$. On the contrary, if $G\cong
K_{k,n-k}$, then $\rho(G)=\sqrt{k(n-k)}.$ This completes the proof. $\Box$
###### Theorem 2.2
Let $G$ be a $P_{t}$-free bipartite graph of order $n$.
(i) If $n<t$, then
$\rho(G)\leq\sqrt{\lfloor\frac{n}{2}\rfloor\lceil\frac{n}{2}\rceil}$. Equality
holds if and only if $G\cong
K_{\lfloor\frac{n}{2}\rfloor,\lceil\frac{n}{2}\rceil}$.
(ii) If $n\geq t$, then $\rho(G)\leq\sqrt{k(n-k)}$, where
$k=\lfloor\frac{t-2}{2}\rfloor$. Equality holds if and only if $G\cong
K_{k,n-k}$ or $G\cong K_{2,2}\cup K_{1}$ for $n=t=5$.
Proof. Since $K_{\lfloor\frac{n}{2}\rfloor,\lceil\frac{n}{2}\rceil}$ has
maximal spectral radius among all bipartite graphs of order $n$, (i) holds
clearly. Next consider (ii).
If $t=2k+2$, that is, $G$ is $P_{2k+2}$-free, then $G$ is also
$C_{2k+2}$-free. Note that $K_{k,n-k}$ does not contain a copy of $P_{2k+2}$.
By (ii) of Theorem 2.1, the result holds.
If $t=2k+3$, then $n\geq 2k+3$. Note that
$\rho(K_{k,n-k})=\sqrt{k(n-k)}\geq\sqrt{k(k+3)}\geq(k+1)=\rho(K_{k+1,k+1}).$
Above inequalities become equalities if and only if $n=2k+3=5$, that is,
$n=t=5$. So we may assume that $G$ does not contain a connected component
isomorphic to $K_{k+1,k+1}$. Then by Lemma 1.2,
$e(N(u),N^{2}(u))\leq(k-1)|N(u)|+k|N^{2}(u)|-k(k-1)$
for any $u\in V(G)$.
Let $X$ be an eigenvector of $G$ corresponding to $\rho(G)$ with
$\sum_{i=1}^{n}x_{i}=1$ and
$B=(b_{ij})_{n\times n}=A^{2}(G)-k(n-k)I.$
Similar to the proof of Theorem 2.1, we have
$\rho^{2}(G)-k(n-k)=\sum_{i=1}^{n}[\rho^{2}(G)-k(n-k)]x_{i}=\sum\limits_{i=1}^{n}\sum\limits_{j=1}^{n}b_{ij}x_{j}=\sum\limits_{u\in
V(G)}b(u)x_{u}\leq 0.$
Hence, $\rho(G)\leq\sqrt{k(n-k)}$. If the equality holds, then $b(u)=0$ for
any $u\in V(G)$. This also implies that $N^{3}(u)=\emptyset$ and (2) becomes
an equality for any $u\in V(G)$. Note that if the case (ii) of Lemma 1.2
occurs, then $d_{G}(u)=k+1$ and $d_{G}(v)=k$. Thus either $G-u\cong
K_{k,|N^{2}(u)|}$ or $G-u\cong K_{|N(u)|,k-1}$ (otherwise, there will exist a
vertex $v\in V(G)$ such that $b(v)<0$). So $G\cong K_{k,n-k}$. we complete the
proof. $\Box$
The following lemma is very useful in characterizing the extremal graph with
minimal least eigenvalue.
###### Lemma 2.3 ([3])
If $G$ is a simple graph with $n$ vertices, then there exists a spanning
bipartite subgraph $H$ of $G$ such that $\lambda_{n}(G)\geq\lambda_{n}(H)$
with equality if and only if $G\cong H$.
Using Lemma 2.3 and Theorems 2.1 and 2.2, we immediately get Theorems 1.3 and
1.4.
Acknowledgment Thanks to Prof. Guanghui Xu who organized a workshop in
Zhejiang A & F University during Aug. 04-10.
## References
* [1] F.K. Bell, D. Cvetković, P. Rowlinson, S.K. Simić, Graph for which the least eigenvalues is minimal, I, Linear Algebra and its Applications, 429 (2008) 234-241.
* [2] F.K. Bell, D. Cvetković, P. Rowlinson, S.K. Simić, Graph for which the least eigenvalues is minimal, II, Linear Algebra and its Applications, 429 (2008) 2168-2179.
* [3] D.S. Cao, Y. Hong, The distribution of eigenvalues of graphs, Linear Algebra and its Applications, 216 (1995) 211-224.
* [4] D. Cvetković, M. Doob, H. Sachs, Spectra of Graphs, third ed., Barth, Heidelberg, 1995.
* [5] Y.Z. Fan, Y. Wang, Y.B. Gao, Minimizing the least eigenvalues of unicyclic graphs with application to spectral spread, Linear Algebra and its Applications, 429 (2-3) (2008) 577-588.
* [6] M. Fiedler, V. Nikiforov, Spectral radius and Hamiltonicity of graphs, Linear Algebra and its Applications, 432 (2010) 2170-2173.
* [7] R.F. Liu, M.Q. Zhai, J.L. Shu, The least eigenvalues of unicyclic graphs with $n$ vertices and $k$ pendant vertices, Linear Algebra and its Applications, 431 (5-7) (2009) 657-665.
* [8] H. Minc, Nonnegative matrices, John Wiley and Sons, New York, 1988.
* [9] V. Nikiforov, A spectral condition for odd cycles in graphs, Linear Algebra and its Applications, 428 (2008) 1492-1498.
* [10] V. Nikiforov, Bounds on graph eigenvalues II, Linear Algebra and its Applications, 427 (2007) 183-189.
* [11] V. Nikiforov, The maximum spectral radius of $C_{4}$-free graphs of given order and size, Linear Algebra and its Applications, 430 (2009) 2898-2905.
* [12] V. Nikiforov, The spectral radius of graphs without paths and cycles of specified length, Linear Algebra and its Applications, 432 (2010) 2243-2256.
* [13] M. Petrović, B. Borovićanin, T. Aleksić, Bicyclic graphs for which the least eigenvalue is minimum, Linear Algebra and its Applications, 430 (4) (2009) 1328-1335.
* [14] Y.Y. Tan, Y.Z. Fan, The vertex (edge) independence number, vertex (edge) cover number and the least eigenvalue of a graph, Linear Algebra and its Applications, 433 (4) (2010) 790-795.
* [15] Y. Wang, Y.Z. Fan, The least eigenvalue of a graph with cut vertices, Linear Algebra and its Applications, 433 (1) (2010) 19-27.
* [16] M.L. Ye, Y.Z. Fan, D. Liang, The least eigenvalue of graphs with given connectivity, Linear Algebra and its Applications, 430 (4) (2009) 1375-1379.
* [17] B.X. Zhu, The least eigenvalue of a graph with a given domination number, Linear Algebra and its Applications, 437 (2012) 2713-2718.
* [18] M.Q. Zhai, B. Wang, Proof of a conjecture on the spectral radius of $C_{4}$-free graphs, Linear Algebra and its Applications, 437 (2012) 1641-1647.
|
arxiv-papers
| 2013-09-26T01:14:21 |
2024-09-04T02:49:51.535567
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Mingqing Zhai, Huiqiu Lin and Shicai Gong",
"submitter": "HuiQiu Lin",
"url": "https://arxiv.org/abs/1309.6700"
}
|
1309.6711
|
# A Spin-Light Polarimeter for Multi-GeV Longitudinally Polarized Electron
Beams
Prajwal Mohanmurthy, and Dipangkar Dutta P. Mohanmurthy and D. Dutta are with
the Department of Physics and Astronomy at Mississippi State University,
Mississippi State, MS 39762.Manuscript received July 11, 2013
###### Abstract
The physics program at the upgraded Jefferson Lab (JLab) and the physics
program envisioned for the proposed electron-ion collider (EIC) include large
efforts to search for interactions beyond the Standard Model (SM) using parity
violation in electroweak interactions. These experiments require precision
electron polarimetry with an uncertainty of $<$ 0.5 %. The spin dependent
Synchrotron radiation (SR), called ”spin-light,” can be used to monitor the
electron beam polarization. In this article we develop a conceptual design for
a “spin-light” polarimeter that can be used at a high intensity, multi-GeV
electron accelerator. We have also built a Geant4 based simulation for a
prototype device and report some of the results from these simulations.
###### Index Terms:
Polarized electrons, synchrotron radiation, spin light, differential
ionization chambers.
## I Introduction
The determination of the longitudinal polarization of the electron beam is one
of the dominant systematic uncertainties in any parity violating electron
scattering (PVES) experiment. In order to achieve the desired high precision,
the polarization of the electron beam must be monitored continuously with an
uncertainty of $<$0.5%. These ambitious goals can be achieved if multiple
independent and high precision polarimeters are used simultaneously. In
addition to being precise, the polarimeters must be non-invasive and must
achieve the desired statistical precision in the shortest time possible.
Compton and Møller polarimeters are typically the polarimeters of choice for
these experiments and are essential to achieve the desired precision. However,
a complimentary polarimetry technique based on the spin dependence of
synchrotron radiation, referred to as “spin-light,” can be used as a relative
polarimeter. A spin-light polarimeter could provide additional means for
improving the systematic uncertainties and when calibrated against a
Compton/Møller polarimeter it could provide a stable continuous monitoring of
the beam polarization. We develop the conceptual design for a continuous
polarimeter based on “spin-light”. The proposed spin-light polarimeter can
achieve statistical precision of $<$ 1% in measurement cycles of less than 10
minutes for 4 - 20 GeV electron beams with a beam currents of $\sim$ 100
$\mu$A.
TABLE I: A comparison of the Compton, Møller and Spin-light polarimeters. Compton | Spin-Light | Møller
---|---|---
non-invasive, | non-invasive, | invasive
continuous | continuous |
analyzing power | analyzing power | analyzing power
energy dependent | energy dependent | energy independent
high currents | moderately | low currents
| high currents |
target is 100% | no target | target is $<$ 10%
polarized | needed | polarized
(requires stable laser) | |
electron & photon | beam left & right | no independent
detection are | detectors provide | measurements
two independent | two independent | possible
measurements | measurements |
high precision | high precision | high precision
absolute polarimeter | relative polarimeter | absolute polarimeter
Best reported [3] | expected | Best reported [1]
instrumental | instrumental | instrumental
uncertainty: 0.4% | uncertainty: 0.6% | uncertainty: 0.47%
Best reported [3] | estimated | Best achieved [2]
absolute | absolute | absolute
uncertainty: 0.5% | uncertainty: $\sim$5% | uncertainty: 0.85%
Møller and Compton polarimeters have a proven track record of very high
precision, the JLab Hall-C Møller polarimeter has an instrumental uncertainty
of 0.47% [1] and absolute uncertainty of 0.85% [2], while the Compton
polarimeter used in the SLD experiment achieved an instrumental uncertainty of
0.4% [3] and an absolute uncertainty of 0.5% [3], hence they are essential for
any PVES program. However, a spin-light polarimeter would have a few
operational and instrumental advantages over conventional polarimeters, such
that when used in parallel with Compton/Møller polarimeters they might help
reduce the systematic uncertainties and achieve the very high precision
essential for the future PVES program. For example, Møller polarimeters use a
polarized Fe target, and the polarization of the target is difficult to
determine and may depend on the beam intensity. Moreover, Møller polarimeters
operate at low current, and are invasive to the primary experiment. Compton
polarimeters require a stable laser (the photon target) and are very sensitive
to backgrounds. The proposed spin-light polarimeter is a target free device,
hence it should be easier to operate over long periods, with its stability
governed just by the stability of the electron beam. Moreover, this novel
polarimeter would facilitate cross-checks and systematic studies when used
with other conventional polarimeters. On the other hand, one of the
disadvantages is that the proposed device can achieve comparable instrumental
uncertainties only as a relative polarimeter, whereas the absolute
polarization is what is required in the PVES experiments. Nevertheless a
precise and stable relative polarimeter can be a very useful device. The spin-
light polarimeter could be used in conjunction with a Compton polarimeter,
such that the difficult to operate Compton polarimeter is used for calibration
and the easier to operate and stable spin-light polarimeter is used to
continuous monitor the beam polarization. Moreover, only the Møller and the
spin-light polarimeters allow measurement of the transverse component of a
longitudinally polarized electron beam. The key features of conventional
polarimeters and a spin-light polarimeter are summarized in Table I.
In this article we discuss the theory behind the spin-light polarimeter and
present a complete conceptual design for such a polarimeter. In order to study
some of the systematic uncertainties of the novel spin-light polarimeter we
have built a Geant4 [4] based simulation of the polarimeter. The results from
these simulations are also presented here. This work was inspired by a 1993
proposal by Karabekov and Rossmanith [5].
## II Synchrotron Radiation
At the typical energies and magnetic fields of present day terrestrial
accelerators, the intensity and angular distribution of SR from an electron
moving along a curved path, under the influence of a magnetic field, is
described with high precision by classical electrodynamics [6]. The total
radiation rate for highly relativistic electrons, $P^{clas}$ is $\propto
E^{4}$ as given by the Larmor formula:
$P^{clas}=\frac{2}{3}\frac{e^{2}\gamma^{4}c}{R^{2}}$, where $e$ is the
electron charge, $v_{e}$ is the velocity of the electron, $c$ is the velocity
of light in vacuum, $\beta=\frac{v_{e}}{c}$,
$\gamma=\frac{1}{\sqrt{1-\beta^{2}}}=\frac{E}{m_{e}c^{2}}$ is the Lorentz
boost, $m_{e}$ is the rest mass of the electron, and, $R$ is the radius of
curvature of the electron orbit. The angular distribution of the radiated
power is given by
$\frac{dP}{d\Omega}=\frac{e^{2}\gamma^{4}c}{4\pi
R^{2}}\frac{(1-\beta\cos\theta)^{2}-(1-\beta^{2})\sin^{2}\theta\cos^{2}\phi}{(1-\beta\cos\theta)^{5}},$
(1)
where the angles $(\theta,\phi)$ are measured with respect to the direction of
the electron’s motion. For highly relativistic electrons ($\gamma>>1$), the
radiation is Lorentz boosted in the forward direction with an opening angle
$\theta\approx 1/\gamma$. The classical theory also tells us that SR is
strongly linearly polarized [7], with $P_{\sigma}=\frac{7}{8}P$ and
$P_{\pi}=\frac{1}{8}P$, where $\sigma$-component is the one where the electric
field lies in the plane of the electron orbit while $\pi$-component is the one
where the electric field lies in the plane perpendicular to the orbital plane.
Experimentally, these properties of SR have been demonstrated with high
precision for a wide range of frequencies [8] \- [11].
Classical theory also tells us that the cone of SR passes over a fixed angle
$\theta$ in retarded 111accounts for the finite delay between the photon
emission and detection time $\Delta
t^{{}^{\prime}}\approx\frac{\Delta\theta}{\omega_{0}}=\frac{1}{\gamma\omega_{0}}$,
where $\omega_{0}=v_{e}/R\approx c/R$ is the Larmor frequency for the
electron. For a distant observer the corresponding time interval is $\Delta
t=\Delta t^{{}^{\prime}}(1-\beta)\approx 1/(2\gamma^{3}\omega_{0})$. Hence the
spectral width of the radiation is $\Delta\omega\approx 1/\Delta
t=2\gamma^{3}\omega_{0}$. This implies that the radiation is strong at high
harmonics of the Larmor frequency $(\omega_{0})$ and can be considered as
continuous. However from quantum mechanics we know that the radiation at
frequency $\omega$ consists of photons of energy $\hbar\omega$. Thus there
must exist a sufficiently strong acceleration such that a single photon will
carry away all of the electron’s energy ($\gamma m_{e}c^{2}=\hbar\omega_{c}$).
The critical magnetic field strong enough to provide this acceleration is
found to be $B_{crit}=\frac{m^{2}c^{3}}{e\hbar}$ = 4.41$\times$109 Tesla. The
frequency of the photon radiated under the influence of $B_{crit}$ is called
the critical frequency ($\omega_{c}$) and is given
by$\omega_{c}=\frac{3}{2}\gamma^{3}\omega_{0}=\frac{3}{2}\gamma^{3}c/R$. From
this we get the critical energy
$E_{crit}=m_{e}c^{2}\sqrt{\frac{m_{e}cR}{\hbar}}\approx$ 106 GeV. The
extremely large values of the critical field and critical energy help explain
why the classical theory is successful at the energies and fields accessible
at present day accelerators. However, it turns out that several quantum
effects appear at considerably lower electron energies.
### II-A Quantum Theory of Synchrotron Radiation
The exact expression for SR intensity including quantum corrections was
calculated by Sokolov, Ternov and Klepikov, based on the solution to the Dirac
equation in the framework of quantum electrodynamics [12]. They showed that at
energies above a few 100 MeV, there would be fluctuations in the radius of the
electron orbit leading to radial oscillations [13] of the electron trajectory
and quantum widening of the trajectory similar to Brownian motion. These
oscillations and widening of the trajectory are essential in determining the
dynamics of electrons in an accelerator, specially storage rings. Sokolov and
Ternov also developed the mathematics required to describe the spin of
relativistic electrons moving in an external electromagnetic field [14, 15],
which allowed them to calculate the electron spin related properties of SR.
The power radiated by electrons undergoing transitions $n\rightarrow
n^{{}^{\prime}}$ (related to the radius of the electron orbit), $s\rightarrow
s^{{}^{\prime}}$ (quadratic fluctuation of the radius) and $j\rightarrow
j^{{}^{\prime}}$ (spin orientations with respect to the magnetic field),
integrated over angles and summed over polarization states, is given by [14,
15]:
$P=P^{clas}\times\frac{9\sqrt{3}}{16\pi}\sum_{s^{{}^{\prime}}}\int_{0}^{\infty}\frac{ydy}{(1+\xi
y)^{4}}I^{2}_{ss^{{}^{\prime}}}(x)F(y),$ (2)
where,
$\displaystyle F(y)$ $\displaystyle=$
$\displaystyle\frac{1+jj^{{}^{\prime}}}{2}[2(1+\xi
y)\int_{y}^{\infty}K_{5/3}(x)dx$ (3) $\displaystyle+$
$\displaystyle\frac{1}{2}\xi^{2}y^{2}K_{2/3}(y)-j(2+\xi y)\xi yK_{1/3}(y)]$
$\displaystyle+$
$\displaystyle\frac{1-jj^{{}^{\prime}}}{2}\xi^{2}y^{2}\left[K_{2/3}(y)+lK_{1/3}(y)\right]$
where $P^{clas}$ is the classical expression for SR power radiated,
$y=\frac{\omega}{\omega_{c}}$, $x=\frac{3}{4}\frac{\xi\gamma^{3}y^{2}}{(1+\xi
y)^{2}}$, $I_{ss^{{}^{\prime}}}$ are the Laguerre functions, $K_{n}(x)$ are
modified Bessel functions, and $\xi=\frac{3}{2}\frac{B}{B_{crit}}\gamma$ is
the critical parameter. Because, $B_{crit}$ = 4.41$\times$ 109 Tesla we have
$\xi<<1$ for magnetic fields used at all man-made accelerators and the above
expression can be expanded in terms of the critical parameter $\xi$ to get
[16];
$\displaystyle P$ $\displaystyle=$ $\displaystyle
P^{clas}[(1-\frac{55\sqrt{3}}{24}\xi+\frac{64}{3}\xi^{2})$ (4)
$\displaystyle-$
$\displaystyle(\frac{1+jj^{{}^{\prime}}}{2})(j\xi+\frac{5}{9}\xi^{2}\frac{245\sqrt{3}}{48}j\xi^{2})$
$\displaystyle+$
$\displaystyle(\frac{1-jj^{{}^{\prime}}}{2})(\frac{4}{3}\xi^{2}+\frac{315\sqrt{3}}{432}j\xi^{2})+...]$
These quantum corrections to the classical expression for charge radiation
($P^{clas}$) involve contributions from the electron recoil effects of
radiation, interference of the charge radiation and radiation due to the
intrinsic magnetic moment of the electron, magnetic moment radiation due to
Larmor precession, magnetic moment radiation due to Thomas precession,
interference between the Larmor and Thomas radiation and radiation due to the
anomalous magnetic moment of an electron [16]. The lowest order spin-dependent
correction is of order $\xi$ and the lowest order spin-flip correction term is
of order $\xi^{2}$. The difference (of order $\xi$) between the expression for
power radiated by polarized and unpolarized (spin averaged) electron beams has
the form,
$\displaystyle P^{pol}-P^{unpol}$ $\displaystyle=$ $\displaystyle-j\xi
P^{clas}\int_{0}^{\infty}\frac{9\sqrt{3}}{8\pi}y^{2}K_{1/3}(y)dy.$ (5)
The above expression is directly related to the spin polarization of an
electron beam $j$, hence the difference $P^{pol}-P^{unpol}=P^{spin}$ can be
called “spin-light” [14, 15]. This offers a new possibility for visual or
direct observation of the polarization characteristics of an electron beam by
determining the SR power at a fixed range of spectral frequency.
The spin dependence of the SR was verified at the VEPP-4 storage ring in
Novosibirsk [17], using a 3 pole wiggler magnet (called a magnetic snake). The
intensity of the SR produced by the wiggler for transversely polarized
electrons was monitored while the beam was periodically depolarized using a RF
field. The measured variation in SR intensity with polarization matched
exactly with the expectation from the Sokolov-Ternov theory. This spin
dependent part of the SR has been successfully used at the VEPP-4 to monitor
the transverse polarization of the electron beam.
The spin-flip term has a unique effect on circular accelerators. The radiation
probability with spin flip is given by
$w^{\uparrow\downarrow}=\frac{1}{\tau}\left(1+j\frac{8\sqrt{3}}{15}\right)$,
where $j=$ 1 is for spin along the magnetic field and $j=$ -1 is for spin
opposite to the magnetic field [14]. Using this relation, Sokolov and Ternov
had predicted that, over time the beam in a circular accelerator would
eventually become polarized opposite to the direction of the magnetic field.
In other words the beam of a circular accelerator becomes transversely
polarized over time via self polarization. This phenomena known as the
Sokolov-Ternov self polarization was first observed at the French storage ring
[18] at Orsay and is now routinely used to polarize beams at circular electron
accelerators such as DESY.
## III Spin Light
In the discussion above the spin orientation $j$ is relative to the magnetic
field that produces the SR, with $j=$ 1 along the magnetic field. Polarized
electron beams have longitudinal ($p_{z}$) and transverse ($p_{\perp}$)
components relative to the beam direction and the transverse polarization have
vertical and horizontal components. For a vertically oriented magnetic field,
the total SR power from transversely polarized electrons, ignoring spin flip
terms and other terms of order $\xi^{2}$, is given by [14]:
$\displaystyle P_{\gamma}(tran)$ $\displaystyle=$
$\displaystyle\frac{9n_{e}}{16\pi^{3}}\frac{ce^{2}}{R^{2}}\gamma^{5}\int_{0}^{\infty}\frac{y^{2}dy}{(1+\xi
y)^{4}}\oint d\Omega(1+\alpha^{2})^{2}$ (6) $\displaystyle\times$
$\displaystyle[K^{2}_{2/3}(z)+\frac{\alpha^{2}}{1+\alpha^{2}}K^{2}_{1/3}(z)$
$\displaystyle-$ $\displaystyle p_{\perp}\xi
y\frac{1}{\sqrt{1+\alpha^{2}}}K_{1/3}(z)K_{2/3}(z)],$
where $n_{e}$ is the number of electrons,
$z=\frac{\omega}{2\omega_{c}}(1+\alpha^{2})^{3/2}$, and $\alpha=\gamma\psi$,
where $\psi$ is the vertical angle in the frame of the moving electron. The
rest of the symbols are as defined in previous instances. The polarization
dependent term in the above expression is an even function of the vertical
angle therefore when integrated over all angles it makes the total SR power
spin dependent. Thus by measuring this spin dependence in the total SR power
radiated one can measure the transverse polarization of the electron beam.
On the other hand, the total SR power from longitudinally polarized electrons,
ignoring spin flip terms and other terms of order $\xi^{2}$, is given by [14]:
$\displaystyle P_{\gamma}(long)$ $\displaystyle=$
$\displaystyle\frac{9n_{e}}{16\pi^{3}}\frac{ce^{2}}{R^{2}}\gamma^{5}\int_{0}^{\infty}\frac{y^{2}dy}{(1+\xi
y)^{4}}\oint d\Omega(1+\alpha^{2})^{2}$ (7) $\displaystyle\times$
$\displaystyle[K^{2}_{2/3}(z)+\frac{\alpha^{2}}{1+\alpha^{2}}K^{2}_{1/3}(z)$
$\displaystyle+$ $\displaystyle p_{z}\xi
y\frac{\alpha}{\sqrt{1+\alpha^{2}}}K_{1/3}(z)K_{2/3}(z)],$
The spin dependent term in the above expression is an odd function of the
vertical angle therefore when integrated over all angles it goes to zero and
the total SR power for longitudinally polarized electrons is spin independent.
However, the power radiated into the space above ($0<\psi<\pi/2$) and below
($-\pi/2<\psi<0$) the orbital plane of the electron are different and the
difference between them is spin dependent. Therefore, by measuring this
spatial asymmetry one can monitor the longitudinal polarization of the
electron beam. For $\xi y<<1$, if we divide the above expression by the energy
of the radiated photon, $E_{\gamma}=\frac{3}{2}\frac{\hbar c}{R}\gamma^{3}y$
we get the total number of photons radiated into a finite horizontal angle
$\Delta\theta$ as,
Figure 1: (a) The total SR power radiated, per MeV (red), and the spin
dependent difference in power radiated, per MeV (blue), above and below the
orbital plane, for 11 GeV, longitudinally polarized electrons in a 4 T
magnetic field and 100 $\mu$A current, and 10% detector efficiency. (b) The
total number of SR photons per MeV and the number of “spin-light” photons per
MeV, and (c) the asymmetry $\frac{\Delta N}{N}$ as a function of photon
energy. (d) The time for 1% statistical uncertainty, in seconds.
$\displaystyle N_{\gamma}$ $\displaystyle=$
$\displaystyle\frac{3}{4\pi^{2}}\frac{1}{137}\frac{I_{e}}{e}\gamma\Delta\theta\int_{y_{1}}^{y_{2}}ydy\int_{-\alpha}^{\alpha}(1+\alpha^{2})^{3/2}$
(8) $\displaystyle\times$
$\displaystyle\left[K^{2}_{2/3}(z)+\frac{\alpha^{2}}{1+\alpha^{2}}K^{2}_{1/3}(z)\right]d\alpha,$
where $I_{e}$ is the beam current. The difference in the photon flux radiated
in the the space above and below the electron orbit is given by,
$\displaystyle\Delta N_{\gamma}(p_{z})$ $\displaystyle=$
$\displaystyle\frac{3}{\pi^{2}}\frac{1}{137}\frac{I_{e}}{e}p_{z}\xi\gamma\Delta\theta\int_{y_{1}}^{y_{2}}y^{2}dy\int_{0}^{\alpha}\alpha(1+\alpha^{2})^{3/2}$
(9) $\displaystyle\times$ $\displaystyle K_{1/3}(z)K_{2/3}(z)d\alpha$
To examine the size and characteristics of the spin dependence we have
numerically integrated the above two expressions for longitudinally polarized
electron with 100% polarization, in a 4 Tesla magnetic field, with $I_{e}$ =
100 $\mu$A, and $E_{e}$ = 11 GeV. We have integrated over a horizontal angular
acceptance of $\Delta\theta$ = 10 mrad, and a vertical acceptance of
$\alpha=\pm$1\. The characteristic spectra of SR and spin-light obtained from
these numerical intragrations are shown in Fig. 1. The total power radiated
$P_{\gamma}$(long) and the spin dependent difference of power radiated above
and below the orbital plane of the electron $\Delta P$(long) are shown as a
function of photon energy in Fig 1 (a). The number of SR photons
$N_{\gamma}$(long) and the number of spin-light photons $\Delta
N_{\gamma}$(long), as function of photon energy are shown in Fig 1 (b). The
asymmetry defined as $A=\frac{\Delta N_{\gamma}(long)}{N_{\gamma}(long)}$ as a
function of photon energy is shown in Fig. 1 (c).
Figure 2: (a) The spin dependent asymmetry as a function of ratio of photon
energy to the critical photon energy, for electron beam energy, Ebeam = 4 – 20
GeV. (b) The spin dependent asymmetry for magnetic field, Bwigg = 2 – 5 Tesla.
Figure 3: (top) The total SR power radiated as a function of the electron beam
energy, for a 100 $\mu$A beam current and Bwigg = 4 Tesla. (bottom) The
vertical size of the SR beam spot as a function of the electron beam energy,
at a distance of 10 m from the wiggler magnet.
Fig. 1 (c) indicates that one should measure the hard tail of the SR spectrum
($E_{\gamma}>$ 500 keV) and avoid the soft part of the spectrum where the
asymmetry is low and changes rapidly with energy. Although the asymmetry is
small $\sim$ 10-4 the photon flux is high, even at the hard tail of the
spectrum, allowing a rapid determination of the asymmetry, with 1% statistical
uncertainty within a few tens of seconds ($\frac{\delta
A}{A}=\frac{1}{A\sqrt{2N}}$) as shown in Fig 1 (d). The energy dependence of
the asymmetry for $E_{e}$ = 4 – 20 GeV and the magnetic field dependence of
the asymmetry for $B_{wigg}$ = 2 – 5 T are shown in Fig 2(a) and Fig 2(b)
respectively. These figures demonstrate that a spin-light based polarimetry is
a very promising technique at intermediate energies and can be used to monitor
the polarization of 4 – 20 GeV electrons in very rapid measurement cycles,
with high statistical precision.
Although the size of the asymmetry increases with increasing electron beam
energy it should be noted that the total power of SR increases as the fourth
power of the electron beam energy and increases linearly with beam current.
Moreover, the vertical size of the SR spot decreases with increasing electron
beam energy as shown in Fig. 3 (bottom). These factors impose practical
limitations on the highest electron beam energies and the highest currents at
which a spin-light polarimeter would be feasible. We estimate that it is best
suited for the 4 - 20 GeV energy range for currents less than $\sim$ 10 mA.
## IV A Conceptual Design
The two basic components of a spin-light based polarimeter are the source of
SR and the X-ray detector which can measure the spatial asymmetry.
### IV-A The SR Source - Wiggler
A three pole wiggler magnet with a magnetic field that has uniform magnitude
but reversed direction at each pole and a short-long-short pole arrangement is
well suited as a source of SR. The three poles must be symmetric about the
center such that the line integral of the magnetic field in the direction of
the motion of the electron, $z$, must be zero (i.e. $\int B(z)dz$ = 0),
ensuring that it does not affect the electron beam transport and its spin
direction (beyond the wiggler). The field being of opposite polarity at the 3
poles, flips the sign of the spin dependent spatial asymmetry from any two
adjacent poles and hence when measured simultaneously it can help reduce
systematic uncertainties arising from the vertical motion of the beam.
The intensity and the asymmetry both increase with increasing field strength,
while the pole length decreases with increasing field strength. Therefore a
field strength of 4 T is a judicious choice for the wiggler field. A 10 mrad
bend can be achieved with a pole length of 10 cm. Thus the total magnet length
is 40 cm, and the spacing between the poles is optimized for ease of
extraction and detection of the SR beam. A separation of 1 m between the poles
allows for collimators to be placed that can separate the SR beams spots from
the different poles. The small pole length ensures that the effect of spin-
flip inducing SR and the fluctuation of the SR power are negligible ($<$
0.1%).
Wiggler magnets are regularly used at light sources around the world such as
the Advanced Photon Source (APS) at Argonne National Lab and Spring8 in Japan.
Some of these magnets are well suited for a spin-light polarimeter [22].
Figure 4: A position sensitive ionization chamber developed at the APS and
SPring-8 with a resolution of 5$\mu$m when operated at photon flux of
5.0$\times$1012 8 keV photons/sec. [23] Note that the image shown here would
constitute just one half of a differential ionization chamber required for
spin-light polarimetry.
#### IV-A1 Effect of the wiggler on the electron beam
A non-invasive polarimeter is highly desirable and hence we must study the
effect of fluctuations related to the quantum nature of SR produced by the
wiggler. The effect of SR on the electron beam were carefully studied for the
recirculating arcs [25], and the same methods can be used to calculate the
influence of the wiggler on critical beam parameters. As described in Ref.
[24] and [25], the distribution of energies lost by individual electrons in
bending through some angle $\theta$ is given by a convolution of the
distribution of the number of photons emitted and the distribution of energies
of those photons. The number of photons emitted by a particular electron per
radian bend will be distributed according to Poisson statistics about a mean
value given by [24], $n=\frac{5}{2\sqrt{3}}\frac{\gamma}{137}=20.62E,$ where
$n$ is the mean number of photons per radian bend, and $E$ is the beam energy
in GeV. The average energy of the photons emitted is
$E_{c}=\hbar\omega_{c}=\frac{3}{2}\frac{\hbar c\gamma^{3}}{R}$ [24]. Therefore
the mean energy fluctuation is given by $\Delta E=\sqrt{n}E_{c}$. It is
interesting to note that the energy fluctuation depends only on the electron
beam energy and the bend radius of the wiggler. A beam of 11 GeV electrons in
a 4T wiggler with a 10 m bend radius and a bend angle of 10 mrad, gives
$n\sim$ 2 and $E_{c}$ = 199 keV. Therefore $\Delta E/E\sim$ 2.5$\times$10-5,
which is comparable to the fluctuations due to the recirculating arcs of the
JLab accelerator [25].
The SR power spectrum usually peaks at angles of $\pm 1/\gamma$ with respect
to the electron direction. However, if an electron emits on the average two
photons in a magnet, the angular distribution of the momentum kick received by
each electron is peaked in the direction of the electron’s motion. The
magnitude of the transverse kicks generated by the emission of a photon with
energy $E_{c}$ in the direction $\theta_{\gamma}=1/\gamma$ with respect to the
electron direction is given by [24],
$\Delta\theta_{e}=\frac{E_{\gamma}\sin{\theta_{\gamma}}}{E_{e}}=11.3\times
10^{-9}\frac{E_{e}(GeV)}{R(m)}$. The r.m.s. kick from the emission of $n$
photons is given by $\sqrt{n}\Delta\theta_{e}$. Thus for a 11 GeV beam bend by
10 mrad the r.m.s. kick is $\sim$ 1.5 $\times$ 10-8 rad, which is negligible.
Thus the wiggler magnet would have negligible influence on the electron beam
and a spin-light polarimeter can be used for non-invasive monitoring of the
beam polarization.
#### IV-A2 Influence of the wiggler bend direction
The wiggler bend direction was chosen to be beam-left for the conceptual
design (see for example Figs. 8 & 9). Since the bend direction is transverse
to the asymmetry direction one does not expect any systematic influence due to
the choice of the bend direction. However, the design includes two symmetric
pairs of ionization chambers placed on either side of the beam, and since the
bend is small (10 mrad), it should be possible to build symmetric pairs of
collimators and slits on both sides of the beam. With such a setup the
independence of the spin-light asymmetry with respect to the wiggler bend
direction can be directly verified during calibration and commissioning of the
device. For the stability of operation, changes in bend direction during
regular operation is not desirable.
### IV-B The X-ray Detector - Ionization Chamber
The detector used to measure the spatial asymmetry must be sensitive to X-rays
in the range of about 500 keV to 2.5 MeV and must be able to pick out a small
asymmetry from a large spin independent background, it must be radiation hard,
have low noise and be able to withstand high rates of $\sim$ 1012 photons/sec.
Ionization chambers (IC) are well known for their high rate capability when
operated as an integrating detector (i.e. in current mode), low electronic
noise and radiation hardness. Argon/Xenon is an attractive candidate for use
as an ionization medium, its high atomic number (18/54) and density (when
compressed) gives it a high stopping power for hard X-rays and low energy
gamma [26]. Over the last two decades, room temperature, high pressure ($>$ 50
atm, 0.55 g/cc) xenon (HPXe) ionization chambers have been developed with high
detection efficiency in the 50 keV - 2.0 MeV range [19, 20, 21].
Figure 5: (left) A schematic of the split anode plate. (right) The dual
differential ionization chamber for spin-light polarimetry.
Another recent development, is the split collector ionization chamber that
have turned the IC into a position sensitive device. Position sensitive
ionization chambers are designed to have the collector plate split into two
sections in a zig-zag/backgammon pattern such that each half operates as an
independent ionization chamber. A prototype of such a chamber has been shown
in Fig. 4. These chambers were developed at the APS at ANL and at the SPring-8
light source in Japan. They are used to measure the vertical position of X-ray
beams and have been shown to have a resolution of 5 $\mu$m [23]. These
chambers also have very low dark currents in the $\sim$pA range and have been
operated at photon flux of 5.0$\times$1012 photons/sec. They work by measuring
the difference in counts between the two halves of the chamber, i.e. they are
differential ionization chambers (DIC). A position sensitive DIC operated in
current mode can be used to measure the spatial asymmetry of the SR generated
by longitudinally polarized electrons.
A dual, 1 atm. Ar/Xe differential ionization chamber would be ideal for a
relative polarimeter. A schematic for such an IC is shown in Fig. 5. The
chamber would consist of Ti or stainless steel windows thick enough to cut
down the low energy X-rays ($<$ 50 keV). A pair of split central anode plates
(separated by a thin insulator) would be placed between the cathodes. The
anode plates would be split in a backgammon pattern. The current measured on
each half of the anode plates is amplified with a differential current
amplifier.
Figure 6: A schematic of the collimated beams of synchrotron radiation from 2
adjacent poles of the wiggler magnet incident on the differential ionization
chambers. The figure on the left is a beam’s view of the electrodes of the
DIC. The collimated radiation is shown as oval blobs, while the figure on the
right is an isometric view of the electrodes, without showing the incident
radiation.
A magnified view of the synchrotron radiation from two adjacent poles of the
wiggler magnet incident on a dual DIC is shown in Fig. 6. The left panel of
the figure shows the beam’s view of the dual DIC. The collimated radiation is
shown as oval blobs with the up-down asymmetry represented by the gray shading
of the blob (the collimation scheme needed to achieve this is discussed in the
next section). The right panel shows an isometric view of the electrodes in
the DIC, without showing the incident radiation. The collimated SR beam from
two adjacent poles will be incident on opposite sides of the anodes in the
dual DIC. The spin-light spatial asymmetry (above and below the orbital plane)
will have opposite sign in each half of the DIC because the magnetic field
direction of the adjacent poles of the wiggler are opposite. On the other hand
any spatial asymmetry due to vertical motion of the beam will have the same
sign in the two halves of the dual DIC and hence should cancel to first order.
Thus the dual DIC is essential to ensure that the spin-light polarimeter is
insensitive to vertical beam motion.
#### IV-B1 The Signal from the DICs
If we denote $N^{L(R)}_{SR}$ as the number of SR photons on the left(right) of
the anode plates, $N^{L(R)}_{spin}$ as the number of spin-light photons, and
$\Delta N^{L(R)}_{z}$ as the difference in number of photons introduce by the
vertical beam motion, then the contribution to the measured current from the
top left part of the dual DIC will be (see Fig. 6);
$I_{1}^{L}\propto N^{L}_{SR}+N^{L}_{spin}+\Delta N^{L}_{z},$
similarly the current contribution from the top right of the dual DIC is;
$I_{1}^{R}\propto N^{R}_{SR}-N^{R}_{spin}+\Delta N^{R}_{z}.$
From the simulation studies we estimate that after collimation the size of
these currents will be on the order of $\sim$ 10 nA.
Note the change in the sign of the contribution from spin light photons
because they are generated from adjacent poles of the wiggler while the
contribution from vertical beam motion has the same sign. For the bottom left
and right parts of the dual DIC we get;
$I_{2}^{L}\propto N^{L}_{SR}-N^{L}_{spin}-\Delta N^{L}_{z},$
and
$I_{2}^{R}\propto N^{R}_{SR}+N^{R}_{spin}-\Delta N^{R}_{z}.$
Thus, the signal from the top and bottom halves of the DIC, $S1$, and $S2$ as
shown in Fig. 6, can be written as,
$\displaystyle S1$ $\displaystyle\propto$
$\displaystyle(N^{L}_{SR}+N^{L}_{spin}+\Delta
N^{L}_{z})-(N^{R}_{SR}-N^{R}_{spin}+\Delta N^{R}_{z})$ (10) $\displaystyle=$
$\displaystyle 2N_{spin},$
and
$\displaystyle S2$ $\displaystyle\propto$
$\displaystyle(N^{L}_{SR}-N^{L}_{spin}-\Delta
N^{L}_{z})-(N^{R}_{SR}+N^{R}_{spin}-\Delta N^{R}_{z})$ (11) $\displaystyle=$
$\displaystyle-2N_{spin}.$
Hence $S1-S2\propto 4N_{spin}\propto 4P_{e}$, and the vertical motion related
asymmetry cancels to first order as does the corrections due to transverse
polarization. However, it should be noted that the signals $S1+S2$ is
proportional to the transverse polarization of the electron beam. This
possibility of measuring both the longitudinal and transverse asymmetries in
the same setup, provides further capability for reducing systematic
uncertainties and makes the spin-light polarimeter an extremely versatile
tool.
Figure 7: (left) The number of photons absorbed in a 1-atm, 50 cm long Xe
chamber. (right) The detector response weighted asymmetry.
The number of photons absorbed in the ionization chamber can be calculated by
multiplying the SR spectrum with the absorption function
$A(\lambda)=1-e^{-\mu(\lambda)\cdot t}$, where $t$ is the length of the
chamber, $\lambda$ is the photon wavelength and $\mu$ is the absorption
coefficient which is obtained from NIST database [27]. The number of photons
absorbed in a 50 cm long chamber with 1 atm Xe, is shown in Fig. 7(left). Also
shown in the absorption weighted (or detector response weighted) asymmetry
(right).
### IV-C Collimation
The spacing between the wiggler poles was chosen to be 1 m to allow adequate
room for the placement of collimators that would separate the SR beams from
each pole. The flight path from the wiggler to the detector is selected to be
10 m, which implies that the SR spot size due to each of the wiggler poles
will fan out over a horizontal length of $\sim$ 10 cm. The vertical width of
the SR spot is only $\sim$ 1 mm. With appropriate placement of collimators on
the wiggler pole entrance and exit faces, it is possible to separate the SR
beam spot from the four different wiggler poles. The magnet system will wiggle
the beam by 10 mrad in the horizontal plane such that each pole of the wiggler
magnet produces a fan of synchrotron radiation in the horizontal plane as
shown in the top view of the magnets (Fig 8).
Figure 8: A schematic of the fan of synchrotron radiation produced as the
electron beam traverses through each of the 4 poles of the wiggler magnet. The
top view of the magnets has been shown. The two colors are used to indicate
that the poles 1 and 4 have opposite polarity compared to poles 2 and 3 and
therefore the sign of the asymmetry for the SR fans of the two colors are
opposite.
For longitudinally polarized electrons this fan of synchrotron radiation will
have an up-down asymmetry in the vertical direction (due to spin-light). A
series of collimators placed at the front and back faces of each pole of the
wiggler magnet and at the center of the central pole of the wiggler will be
used to select small angular ranges from the entire fan of synchrotron
radiation as shown in the top view of the magnets in Fig 9. This collimation
scheme enables separation of the synchrotron radiation from each pole of the
wiggler magnet. Such a separation is necessary because the up-down asymmetry
of the synchrotron radiation has opposite sign for each pole of the wiggler.
Each of the collimated beams of synchrotron radiation will be separated by a
few cm when they are projected into two symmetric dual DICs located at a
distance of 10 m from the wiggler magnet.
Figure 9: A schematic for the slits and collimators used to select small
angular range of the fan of synchrotron radiation. The synchrotron beams from
poles 1 and 2 will be incident on the dual DIC placed on the beam left and the
beams from poles 3 and 4 are incident on the dual DIC placed on the beam
right.
### IV-D The Complete Polarimeter
Figure 10: A schematic for a differential spin-light polarimeter (not to
scale).
A 3D view of the complete spin-light polarimeter is shown in Fig 10. Each pole
of the wiggler magnet is separated by a distance of 1 m and the two dual
differential ionization chambers are placed 10 m from the last pole of wiggler
magnet. The vertical backgammon split on the central anodes makes the DIC
position sensitive in the vertical direction and hence the signal from the DIC
is sensitive to the up-down asymmetry for each of the collimated beams of
synchrotron radiation. The second dual DIC is necessary to provide an
independent measurement of the up-down asymmetry and help reduced systematic
uncertainties. The main parameters of this conceptual design are tabulated in
Table II.
TABLE II: Parameters of the SR polarimeter at 11 GeV | Beam & Magnet
---|---
Ee, current | (GeV), ($\mu$A), | 11.0, 100
Bwiggler | (T) | 4.0
Pole (total) length | (m) | 0.1 (0.4)
Separation | (m) | 1.0
between poles | |
bend angle | (mrad) | 10
vert. opening | (mrad) | 0.05
angle | |
Flight path | (m) | 10
to det. | |
| SR and detector
Nγ/s | (Hz) | 5.8$\times$ 1013
$\Delta$NSpin/s | (Hz) | 1.8$\times$ 109
Detector | (1 atm Xe, cm) | 50
medium | |
Nabs/s | (Hz) | 3.1$\times$ 1012
$\Delta$Nfluctuation/s | (Hz) | 7.6$\times$ 106
vert. beam | (mm) | 1.0
spot after 10m | |
flight | |
### IV-E Systematic Instrumental Uncertainties
Some of the major sources of systematic uncertainties for a spin-light
polarimeter include the background asymmetries from processes such as
Bremsstrahlung and false asymmetry due to vertical beam motion, differences in
chamber efficiency and magnetic field non-uniformity between adjacent poles of
the wiggler. The measured experimental asymmetry from a spin light polarimeter
can be written as $A_{expt}=A_{raw}(1+B/S)-A_{B}B/S+A_{F}$, where $S$ and $B$
are the signal and background, $A_{B}$ is the background asymmetry and $A_{F}$
is the false asymmetry due to factors such as vertical beam motion,
differences in the chamber efficiency and differences in the field strength
between adjacent poles. The main advantage of operating the ionization
chambers as differential detectors is that the false asymmetries will cancel
to first order. In addition the visible portion of the synchrotron light can
be used to align the detectors and help control systematic uncertainties. The
3-pole design ensures that the vertical beam motion related false asymmetry
also cancels to first order. However, the size of the background must be small
compared to the signal. In order to address this issue, a full Geant4 [4]
simulation of a proto-type spin-light polarimeter was built. The simulation
was also used to study the effect of the asymmetry associated with the
background. In addition to the built-in synchrotron radiation physics
available in Geant4, we have implemented a spin-light generator using a
parametrized model of spin-light at the nominal running conditions (11 GeV
beam, 4T magnetic field and 10 mrad bend angle). Using this generator we were
able to reproduce the expected photon energy spectrum and the expected
asymmetry as shown in Fig. 11.
Figure 11: The simulated photon spectrum and asymmetry compared to the
calculated spectra for the nominal running conditions.
The exact position and width of the slits/collimators on the front and back
faces of the wiggler magnets were optimized iteratively to obtain the best
separation between the SR photons originating from the different poles of the
magnets. Currently, we are using slits only to illuminate the beam right split
ionization chamber. A 3 mm sheet of lead shielding is also applied in front of
the ion chamber.
The ionization generated from the photons incident on the ion-chamber is
integrated to obtain the signal from the chamber and this signal was used to
study several of the systematic uncertainties, such as the effect of
background radiation, the position/alignment of the collimators and the
alignment of the ion chambers, on the asymmetry signal. The results of these
studies have been tabulated in Table III.
The simulated data show that the background from non-SR radiation is $\sim$
1.6% and most of the background is at energies below 0.5 MeV. If the
background related dilution can be determined to $\sim$ 1% the systematic
uncertainty due of background radiation will be $\sim$ 0.5%. Moreover, in the
proposed setup the background can be determined by measuring the difference in
the signal from the chambers with the wiggler magnets turned on and off during
calibration and commissioning of the device and at some relatively long
interval during regular operations (these background measurements would be
invasive in nature).
The same Geant4 simulation was also used to determine the contribution to the
systematic uncertainty due to the uncertainty in the exact position and width
of the slits and collimators that are placed on the front and back faces of
the poles of the wiggler magnet. These collimators are used to separate the SR
photons originating from the various poles. Thus, any uncertainty in the
position and width of the collimators can give rise to a false asymmetry due
to mixing of spin-light and SR photons from different poles. The simulation
demonstrates that a position/width uncertainty of 100 $\mu$m would lead to a
uncertainty of $<$ 0.2%.
We have also simulated the effects of the finite beam size on the up-down
asymmetry. The simulations were performed for a beam with a Gaussian
distribution of $\sigma$= 100 $\mu$m. The effects of the fringe fields at the
tapered edges of the wiggler poles has also been studied using a Poisson Fish
[28] model of the magnetic field. Although the absolute photon flux was
reduced due the fringe fields the effect on the asymmetry was found to be
minimal.
We have studied the effects of beam halos on the spin-light asymmetry. The
restrictive collimation scheme discussed in section IV-C ensures that the
contribution from beam halo is very limited. For a ratio of peak to halo of
$\sim~{}10^{8}$ the dilution to the spin-light asymmetry is several orders of
magnitude smaller than the spin-light asymmetry. Thus, for a tightly
controlled beam, halos should not pose a serious challenge.
Most of the systematic uncertainties listed in Table III are independent of
the electron beam energy in the 4 - 20 GeV range, however, because of the
decrease in the SR spot size with increasing electron beam energy the
uncertainty due the slit width scales inversely with beam energy.
The proposed device is best used as a relative polarimeter, however, it can be
used as an absolute device if the lower and upper bounds of the energy
sensitivity of the DIC is determined accurately. The absolute value of the
spin light asymmetry depends on the absolute value of the energy window over
which the DIC signals are integrated. It is especially sensitive to the lower
bound because of the rapid change in the SR intensity and the spin-light
asymmetry as a function of decreasing photon energy. For an absolute
measurement the lower bound of the integration window and the sensitivity to
the vertical motion of the beam would be the two dominant sources of
systematic uncertainty. Excellent energy resolutions have been demonstrated
for HPXe ionization chambers [19, 20, 21]. With such high resolution
ionization chambers one should be able to determine the response function and
the lower bound of energy sensitivity of the chamber to better than 2%. Using
a lower bound of 0.5 MeV, a variation of $\pm$ 10 keV in the lower bound
results in a 2.5% change in the calculated asymmetry. An absolute measurement
would be sensitive to vertical beam motion, a $\pm$ 0.1 mm variation in the
vertical beam position can result in a 4% change in the calculated asymmetry.
Thus, a spin-light polarimeter would only be capable of $\sim$ 5% absolute
polarization measurement.
A table of estimated systematic uncertainties is shown in Table III. We
estimate the systematic instrumental uncertainties of a relative polarimeter
to be $<$ 1%.
TABLE III: Systematic instrumental uncertainties for a relative polarimeter at $E_{beam}=11$ GeV. Source | Uncertainty | $\frac{\delta A}{A}$
---|---|---
Dark current | $\sim$ pA | $<$ 0.01%
Intensity fluctuations | $\Delta N\times$ 10-3 | $<$0.1 %
Beam energy | 1.0$\times$10-3 | $<$ 0.05 %
Density of chamber gas, | relative difference | $<$0.01%
Slit width | 100 $\mu$m | $<$0.2 %
Background related | known to 0.5% | 0.5 %
dilutions | for B/S $\sim$ 0.02 |
Other dilutions | cancel to first order | $<$ 0.1%
Halo contributions | $10^{-8}$ | $<$ 0.1 %
Total | | 0.6 %
## V Conclusion
Spin light based polarimetry was demonstrated over 30 years ago, but has been
ignored since then. A spin-light polarimeter has several advantages over
conventional polarimeters and when used in conjunction with a Compton
polarimeter it could help provide a new benchmark for precision polarimetry.
The 11 GeV beam at JLab or the electron beam at a future EIC would be well
suited for spin light polarimetry and such a polarimeter would help achieve
the $<$ 0.5 % polarimetry desired by experiments approved for the 12 GeV era
and proposed for the EIC. A 3 pole wiggler with a field strength of 4 T and a
pole length of 10 cm would be adequate for such a polarimeter. A dual position
sensitive ionization chambers with split anode plates is ideally suited as the
X-ray detector for such a polarimeter. The differential detector design would
help reduce systematic uncertainties.
## Acknowledgment
This work was supported in part by the U.S. Department of Energy under
contract # DE-FG02-07ER41528, by the EIC Detector R&D grant from Brookhaven
National Lab. One of us (P.M.) would also like the thank the Jefferson Science
Associates for a JSA Fellowship.
## References
* [1] M. Hauger et al., Nucl. Inst. Meth. A462, 382 (2001).
* [2] To be reported for the polarization measured during the QWeak experiment at JLab, Hall-C.
* [3] P. C. Rowson, D. Su, and S. Willocq, Ann. Rev. Nucl. Part. Sci. 51, 345 (2001); M. Woods, SLAC-PUB-7319 (1996).
* [4] J. Allison et al., Nucl. Inst. and Meth. A506, 250 (2003); J. Allison et al., IEEE Trans. in Nucl. Science 53, 270 (2006).
* [5] I. P. Karabekov, R. Rossmanith, Proc. of the 1993 PAC, Washington, v. 1, p. 457 (1993); I. P. Karabekov and S. I. Karabekian, Proceedings of 5th European Particle Accelerator Conference (EPAC 96), Sitges, Spain, 10-14 Jun 1996, pp 1743-1745 (1996); A. V. Airapetian, R. O. Avakian, I. P. Karabekov, E. L. Saldin, and M. V. Yurkov, Proc. of the SPIN-96, Amsterdam, The Netherlands, Vol.1, p762 (1996).
* [6] D. D. Ivanenko, I. Pomeranchuk, Ya Zh. Eksp. Teor. Fiz. 16, 370 (1946); J. Schwinger, Phys. Rev. 75, 1912 (1947).
* [7] G. A. Schott, Ann. Phys. 24, 635 (1907); A. A. Sokolov and I. M. Ternov, Zh. Eksp. Theor. Fiz. 31, 373 (1956), Sov. Phys. JETP 4, 396 (1957).
* [8] F. R. Elder, R. V. Langmuir and H. C. Pollock, Phys. Rev. 74, 52 (1948).
* [9] M. Yu Ado and P. A. Cherenkov, Sov. Phys. Dokl. 1, 517 (1957).
* [10] F. A. Korolev, E. .N. Akimov, E. N. Markov, and O. F. Kulikov Sov. Phys. Dokl. 1, 568 (1957).
* [11] P. Joos, Phys. Rev. Lett. 4, 558 (1960).
* [12] A. A. Sokolov, N. P. Klepikov and I. M. Ternov, JETF 23, 632 (1952).
* [13] A. A. Sokolov, and I. .M. Ternov, JETF 25, 698 (1953).
* [14] A. A. Sokolov and I. .M. Ternov, Synchrotron Radiation, Pergamon Press, New York (1968); A. A. Sokolov and I. .M. Ternov, Radiation from Relativistic Electrons, A.I.P. Translation Series, New York (1986).
* [15] I. M. Ternov, Physics - Uspekhi 38, 409 (1995).
* [16] V. A. Bordovitsyn, Ph. D. Thesis, Moscow (1983); I. M. Ternov and V. A. Bordovitsyn, Vestn. Mosk. Univ. Ser. Fiz. Astr. 24, 69 (1983); V. A. Bordovitsyn and V. V. Telushkin, Nucl. Inst. and Meth. B266, 3708 (2008).
* [17] S. A. Belomesthnykh et al., Nucl. Inst. and Meth. 227, 173 (1984).
* [18] J. Le Duff, P. C. Marin, J. L. Manson, and M. Sommev, Orsay - Rapport Technique, 4-73 (1973).
* [19] V. V. Dmitrenko et al., Sov. Phys.-tech. Phys. 28, 1440 (1983); A. E. Bolotnikov et al., Sov. Phys.-Tech. Phys. 33, 449 (1988);
* [20] C. Levin et al., Nucl. Inst. and Meth. A332, 206 (1993).
* [21] G. Tepper and J. Losee, Nucl. Inst. and Meth. A356, 339 (1995).
* [22] E. Nakamura et al., J. of Elec. Spec. and Rel. Phen. 80, 421 (1996); D. E. Baynham, P. T. M. Clee, and D. J. Thompson, Nucl. Instr. and Meth., 152, 31 (1978)
* [23] K. Sato, J. of Synchrotron Rad., 8, 378 (2001); T. Gog, D. M. Casa and I. Kuzmenko, CMC-CAT technical report.
* [24] M. Sands, SLAC Technical note, SLAC-121 (1970).
* [25] B. Norum, CEBAF Technical note, TN-0019 (1985).
* [26] A. E. Bolotnikov and B. Ramsey, Nucl. Inst. and Meth. A396, 360 (1997).
* [27] http://www.nist.gov/physlab/data/xraycoef/index.cfm
* [28] http://library.lanl.gov/cgi-bin/getfile?00415886.pdf
Prajwal Mohanmurthy Prajwal Mohanmurthy obtained his Bachelor of Science
degree from Mississippi State University in 2012. He was a graduate research
fellow at the High Performance Computing Collaboratory at Mississippi State
University and currently is a graduate research fellow in the Laboratory for
Nuclear Sciences at Massachusetts Institute of Technology. His research
interests are centered around test of standard model and fundamental
symmetries in search of physics beyond the standard model. His recent research
involvements have been geared towards a search for Axionic Dark Matter and the
precision measurement of the mass of neutrinos. He also actively collaborates
to develop beam instrumentation for up and coming facilities and future
accelerators.
---
Dipangkar Dutta Dr. Dipangkar Dutta is an Associate Professor of Physics at
the Mississippi State University Department of Physics and Astronomy. He
obtained his Bachelor of Technology degree from Indian Institute of
Technology, Bombay in 1992 and his doctoral degree in Physics from
Northwestern University in 1999. He was a post-doctoral and senior post-
doctoral fellow in the Laboratory for Nuclear Sciences at Massachusetts
Institute of Technology. His research is focused primarily on precision
measurement of fundamental properties of nucleons. He is also interested in
precision tests of fundamental symmetries and the Standard Model.
---
|
arxiv-papers
| 2013-09-26T02:22:48 |
2024-09-04T02:49:51.542174
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Prajwal Mohanmurthy, Dipangkar Dutta",
"submitter": "Prajwal Mohanmurthy",
"url": "https://arxiv.org/abs/1309.6711"
}
|
1309.6773
|
# PyECLOUD and build-up simulations at CERN
G. Iadarola(1),(2) and G. Rumolo(1)
(1) CERN Geneva Switzerland (2) Università di Napoli “Federico II Naples
Italy
###### Abstract
PyECLOUD is a newly developed code for the simulation of the electron cloud
(EC) build-up in particle accelerators. Almost entirely written in Python, it
is mostly based on the physical models already used in the ECLOUD code but,
thanks to the implementation of new optimized algorithms, it exhibits a
significantly improved performance in accuracy, speed, reliability and
flexibility. Such new features of PyECLOUD have been already broadly exploited
to study EC observations in the Large Hadron Collider (LHC) and its injector
chain as well as for the extrapolation to high luminosity upgrade scenarios.
## 1 INTRODUCTION
The analysis of the electron cloud (EC) observations in the Large Hadron
Collider (LHC) and its injectors have raised new challenges for the EC build-
up simulations. On one hand, for a correct understanding of machine
observations it is often necessary to deal with beams with thousands of
bunches and with non idealities like non uniform bunch populations and bunch
lengths along the train. On the other hand, the demand for extensive
parametric scans gives quite stringent requirements in terms of speed and
reliability.
CERN has a long experience in the EC build-up simulation, mostly carried out
with the ECLOUD code, developed and maintained at CERN since 1997 [1, 2, 3,
4]. Unfortunately, due to its not modular structure and to the programming
language (FORTRAN 77), this code would need a deep reorganization and a
serious upgrade to enable it to fulfill the aforementioned requirements.
Therefore we have decided to write a fully reorganized code, in a more modern
and powerful language, considering that the initial effort would be
compensated by a significantly increased efficiency in future developments and
debugging. The new code has been called PyECLOUD, since it is almost entirely
written in Python and is largely based on the physical models of the ECLOUD
code. Nevertheless, several features and implementations have been modified,
in some cases completely redesigned, with respect to ECLOUD, with substantial
improvements in terms of reliability, accuracy, speed and usage flexibility.
## 2 PyECLOUD
Like ECLOUD, PyECLOUD is a 2D macroparticle (MP) code, where the electrons are
grouped in MPs in order to achieve a reasonable computational burden.
The dynamics of the MP system is simulated following the flow diagram sketched
in Fig. 1.
Figure 1: Flowchart representing PyECLOUD main loop. Figure 2: Top: evolution
of the number of electrons in the beam pipe for an LHC type beam with 25ns
bunch spacing in the SPS (2 trains of 72 bunches); middle: evolution of the
reference MP size; bottom: evolution of the number of MPs, the regeneration
threshold is highlighted in red.
At each time step, seed electrons, due to residual gas ionization and/or to
synchrotron radiation induced photoemission from the chamber walls, are
generated with a number consistent with the passing beam slice and with
positions and momenta determined by theoretical or empirical models.
Then the electric field acting on each MP is evaluated: the field of the beam
is precomputed on a suitable rectangular grid, loaded from file and obtained
at each MP location by a linear (4 points) interpolation; the space charge
contribution of the electron system itself is calculated by a classical
Particle in Cell (PIC) algorithm, where the finite difference method is
employed to solve the Poisson equation with perfectly conducting boundary
conditions on the beam chamber.
Once the total electric field at each MP location is known, MP positions and
momenta are updated by integrating the dynamics equation; at this stage the
presence of an externally applied dipolar magnetic field can also be taken
into account.
At each time step, a certain number of MPs can hit the wall. In these cases a
proper model of the secondary emission process is applied to generate charge,
energy and angle of the emitted electrons. According to the size of the
emitted charge, a rescaling of the impinging MP can be performed or new MPs
can be emitted.
### 2.1 MP size management
One of the peculiarities of the EC build-up process is the fact that, due to
the multipacting effect, the number of electrons can spread several orders of
magnitude during the passage of the bunch train (see Fig. 2 \- top). As a
consequence, it is impossible to choose a MP size which is suitable for the
entire simulation, allowing both a satisfactory description of the phenomena
and a computationally affordable number of MPs. The MP size management in
PyECLOUD has been significantly improved with respect to ECLOUD and will be
briefly described in this subsection.
MP sizes are not enforced throughout the simulation process but are determined
step by step by “decisions” taken during the execution. For this purpose a
target MP size $N_{ref}$, dynamically adapted during the simulation, is
employed to control the number of electrons per MP. In particular:
* •
The size of MPs generated by seed mechanisms is exactly $N_{ref}$;
* •
When a MP hits the wall, it is simply rescaled according to the Secondary
Electron Yield (SEY) if the emitted charge is below $1.5N_{ref}$, otherwise
”true” secondary MPs are generated so that the resulting MP size is as close
as possible to $N_{ref}$;
* •
Once per bunch passage, a cleaning procedure is performed, which deletes the
MPs with charge lower than $10^{-4}N_{ref}$.
$N_{ref}$ is changed whenever the total number of MPs becomes larger than a
certain threshold defined in the input file(typical value $\sim 10^{5}$),
which means that the computational burden has become too high. When this
happens, a regeneration of the set of MPs is applied, by the following
procedure (see Fig. 2):
* •
Each MP is assigned to a cell of a uniform grid in the 5-D phase space
$(x,y,v_{x},v_{y},v_{z})$ obtaining the phase space distribution of the
electron gas distributed on the mesh points;
* •
The new $N_{ref}$ is chosen in order to get a target number of MPs (typically
5-10 times smaller than the regeneration threshold), which still allows for an
accurate simulation but with a more reasonable computational effort;
* •
A new set of MPs, having the new reference size, is generated according to the
computed distribution.
The preservation of the entire phase space is very important in EC build-up
simulation since the dynamics imparted by passing bunches generates very
distinctive velocity distributions at the different time steps and the
conservation of few specific moments would not guarantee a sufficient
accuracy.
Several numerical test have shown that the errors on the total charge and the
total energy which are introduced by this procedure, are about 1% at the first
time step after the regeneration and they become even smaller at the first
bunch passage after the regeneration.
Figure 3: Electron cloud build-up simulations for different time steps in ECLOUD (top) and PyECLOUD (bottom). Simulated case: SPS MBB bending magnet, 26GeV, two trains of 72 bunches with 225ns gap, 25ns bunch spacing, 1.1$\cdot 10^{11}$ protons per bunch. | Processing time
---|---
Time step [ps] | ECLOUD | PyECLOUD
200 | 29 min | 12 min
100 | 1 h 27 min | 13 min
50 | 1 h 45 min | 24 min
25 | 3 h 7 min | 40 min
12 | 4 h 15 min | 1 h 6 min
Table 1: Computation time required by ECLOUD and PyECLOUD for the simualtions
in Fig. 3.
### 2.2 Performances
The passage from ECLOUD to PyECLOUD had a significant impact on the
performances both in terms of accuracy and of computational efficiency. Fig. 3
shows a comparison between the two codes in terms of convergence properties
with respect to the time step which is chosen for the simulation. While in
ECLOUD it is quite difficult to achieve a good convergence, PyECLOUD gives a
good estimate of the total number of electrons in the chamber already for
quite large time steps ($\sim 0.1$ns) while a satisfactory convergence is
obtained for a time step of the order of $25$ps.
For the same test cases the simulation time required by the two codes is
reported in Tab. 1 showing that the improvements introduced in PyECLOUD had
also a positive impact on the code efficiency.
Furthermore, the new code has been designed in order to offer an increased
usage flexibility, allowing to deal with irregular beam structures e.g. non
uniform bunch intensity and/or bunch length along the bunch train, irregular
bunch spacings and bunch profiles.
Thanks to these new features, PyECLOUD has been already largely exploited at
CERN for several EC simulation studies for the LHC and its injector chain [5,
6, 7, 8]. In particular, as described in detail in [9], PyECLOUD simulations
have been used to reconstruct the evolution of the SEY of the chambers in the
LHC arcs, from the measurement of the heat load deposited on the beam screen
of the cryogenic magnets. The new code also allows us to estimate the bunch by
bunch energy loss due to the interaction of the beam with the EC and to export
the electron distribution seen by each bunch. The first feature allowed us to
benchmark the results against bunch by bunch stable phase measurements [10]
while the second was used, together with HEADTAIL simulations, to analyze the
instabilities observed in the LHC with 25ns bunch spacing [11].
Simulation studies have also addressed the EC formation in the common vacuum
chambers of the LHC. An example of this kind of application will be described
in the following section in order to show the capability of PyECLOUD to deal
with beams made of thousands of bunches with irregular spacings.
## 3 EC build up in LHC common vacuum chambers
Common vacuum chambers having 800mm diameter are installed on both sides of
the ALICE experiments in the Long Straight Section 2 (LSS2) of the LHC. During
2011 operation with 50ns bunch spacing an important pressure rise was noticed
in these chambers with a significant impact on the background observed by the
ALICE experiment.
The analysis of the pressure data has shown that a severe pressure increase is
observed only when the two rings of the LHC are completely filled. Fig. 4
shows the pressure evolution during a proton-proton physics fill (with 50ns
spacing) in which the injection from the SPS of the last two trains of 144
bunches was delayed by about one hour. It can be noticed that the pressure
rise appears already at the injection energy (450GeV), but only after the last
two injections have taken place.
Figure 4: Pressure evolution in the 800mm chambers near the ALICE experiment
of the LHC during a proton physics fill. The total beam intensity is the two
rings is also included.
PyECLOUD simulations have been run in order to investigate if the EC formation
in these chambers could explain this pecular behavior. The presence of both
counter-rotating beams in the chamber had to be taken into account since it
determines different “hybrid” bunch spacings at the different sections of the
$\sim$30m long vacuum chambers when both beams are circulating in the LHC (see
Fig. 5).
In particular the two beam configurations in Fig. 6 have been simulated, which
correspond to the beam patterns in the two rings at the moments indicated by
(a) and (b) in Fig. 4. The results at a certain section of the considered
vacuum pipes are shown in Fig. 7.
It can be noticed that, in the configuration of Fig. 6a, both beams present a
gap of about one quarter of the length of the ring. Probably due to the quite
large radius, the EC can develop only when both beams are passing in the
chamber, while a decay of the number of electrons is observed when only one
beam is passing and the gap is long enough to allow a complete reset of the EC
between subsequent turns.
On the other hand, after the injection of the last two trains from the SPS,
the layout of the beam in the two rings looks like the one in Fig. 6b where no
large gap is present in any of the two beams. As a consequence a complete
decay of the EC between subsequent turns is not possible anymore. In fact, a
memory effect is observed between turns with a strong enhancement of the EC
activity and, as consequence, of the electron stimulated gas desorption
leading to the observed pressure rise.
Figure 5: Hybrid bunch spacings which are observed at different sections of
the 800mm common chambers. Figure 6: Filling pattern of the two beams in the
LHC before (a) and after (b) the injections of the last two trains. Beam 1
(blue) is clockwise rotating, beam 2 (red) counterclockwise. The position the
straight section 2 (ALICE) is highlighted Figure 7: PyECLOUD simulation of the
electron cloud build-up in the 800mm common chambers for the filling patterns
in Fig. 6.
## 4 CONCLUSIONS
PyECLOUD is a new EC build-up simulation code, which has been developed at
CERN for a reliable and efficient analysis of the EC observations in LHC as
well as for EC studies related to the high luminosity upgrade of the LHC and
its injectors.
Based on the physical models already implemented in ECLOUD, the new code
features several improvements in terms of algorithms and implementation (e.g.
a new MP size management method) with a significant impact on speed and
convergence properties. Moreover the code is explicitly designed to deal with
irregular beam structures (e.g. different bunch intensities and bunch lengths
along the train, arbitrary spacings and profiles) in order to allow for an
accurate analysis of EC observations in CERN accelerators.
Several EC studies for the LHC and its injectors have been already carried out
with the new code giving very encouraging indications on the reliability of
the models and numerical solutions.
## 5 ACKNOWLEDGMENTS
The authors would like to express their gratefulness to G. Arduini, H.
Bartosik, C. Bhat, V. Baglin, R. De Maria, O. Dominguez, M. Driss Mensi, J.
Esteban-Muller, K. Li, H. Maury Cuna, G. Miano, E. Méétral, H. Neupert, G.
Papotti, E. Shaposhnikova, M. Taborelli, L. Tavian, C. Y. Vallgren, and F.
Zimmermann for the support they provided in the code development, the
simulation work and the machine data analysis as well as for valuable comments
and discussions.
## References
* [1] F. Zimmermann, “A Simulation Study of Electron-Cloud Instability and Beam-Induced Multipacting in the LHC”, CERN LHC Project Report 95 (February 1997)
* [2] F. Zimmermann, “Electron Cloud Simulations for SPS and LHC”, Proc. Chamonix X, 17-21 January 2000, CERN-SL-2000-007 DI (2000), p. 136
* [3] G. Rumolo and F. Zimmermann, “Practical User Guide for ECloud”, CERN-SL-Note-2002-016.
* [4] G. Rumolo et al.,“Simulation of the electron-cloud build up and its consequences on heat load, beam stability, and diagnostics”, Phys. Rev. ST Accel. Beams 4 (2001), 1
* [5] G. Iadarola and G. Rumolo, “Electron Cloud in the CERN accelerators (PS, SPS, LHC)”, elsewhere in these proceedings
* [6] H. Maury Cuna et al., “Simulation of electron-cloud heat load for the cold arcs of the Large Hadron Collider”, Proceedings of the 3rd International Particle Accelerator Conference 2012, 20 - 25 May 2012, New Orleans, USA
* [7] F. Caspers et al., “Comparison between electron cloud build-up measurements and simulations at the CERN PS”, Proceedings of the 3rd International Particle Accelerator Conference 2012, 20 - 25 May 2012, New Orleans, USA
* [8] C. Bhat et al., “E-Cloud dependence on the Bunch Profile - An Experiment in the PS and an Extension to the LHC”, elsewhere in these proceedings
* [9] G. Iadarola and G. Rumolo, “Electron cloud simulations with PyECLOUD”, Proceedings of the 11th International Computational Accelerator Physics Conference (ICAP), 19-24 August 2012, Rostock-Warnemünde, Germany
* [10] J. Esteban Muller et al., “Synchronous Phase Shift at LHC”, elsewhere in these proceedings
* [11] H. Bartosik et al., “Benchmarking HEADTAIL with e-cloud observations with LHC 25ns beam”, elsewhere in these proceedings
|
arxiv-papers
| 2013-09-26T09:41:16 |
2024-09-04T02:49:51.552996
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "G. Iadarola (CERN and U. Naples), G. Rumolo (CERN)",
"submitter": "Scientific Information Service CERN",
"url": "https://arxiv.org/abs/1309.6773"
}
|
1309.6795
|
# ELECTRON CLOUD IN THE CERN ACCELERATORS (PS, SPS, LHC)
G. Iadarola(1+2) and G. Rumolo(1)
(1) CERN Geneva Switzerland (2) Università di Napoli “Federico II Naples
Italy
###### Abstract
Several indicators have pointed to the presence of an Electron Cloud (EC) in
some of the CERN accelerators, when operating with closely spaced bunched
beams. In particular, spurious signals on the pick ups used for beam
detection, pressure rise and beam instabilities were observed at the Proton
Synchrotron (PS) during the last stage of preparation of the beams for the
Large Hadron Collider (LHC), as well as at the Super Proton Synchrotron (SPS).
Since the LHC has started operation in 2009, typical electron cloud phenomena
have appeared also in this machine, when running with trains of closely packed
bunches (i.e. with spacings below 150ns). Beside the above mentioned
indicators, other typical signatures were seen in this machine (due to its
operation mode and/or more refined detection possibilities), like heat load in
the cold dipoles, bunch dependent emittance growth and degraded lifetime in
store and bunch-by-bunch stable phase shift to compensate for the energy loss
due to the electron cloud.
An overview of the electron cloud status in the different CERN machines (PS,
SPS, LHC) will be presented in this paper, with a special emphasis on the
dangers for future operation with more intense beams and the necessary
countermeasures to mitigate or suppress the effect.
## 1 Introduction
In the CERN PS, the electron cloud was first observed in 2001 during the last
part of the cycle for the production of the the so-called LHC-type beams, i.e.
the beams of the type needed for the LHC filling. The production scheme of
these beams in the PS is based on two or three steps of bunch splitting in
order to obtain at the exit of the PS bunch trains with 50ns or 25ns spacing,
respectively. In either case, the final stage of bunch splitting takes place
at the top energy (26 GeV/c) and is followed by adiabatic bunch shortening and
fast bunch rotation shortly before extraction [1]. These two processes are
meant to shorten the bunches from their 15 ns length after the last splitting
to 12 and then 4 ns, respectively, and make them suitable to be injected into
the SPS. Therefore, these beams only circulate in the PS for few tens of msec
with a structure prone to electron cloud formation (beam parameters are
summarized in Table 1).
During this short time before extraction, an electron cloud was initially
revealed in 2001 by the presence of a baseline drift in the signal from the
pick up as well as beam transverse instabilities [2]. The transverse
instabilities made a new appearance with 25ns beams in 2006, when the bunches
were accidentally shortened to 10ns or below (instead of the nominal 12ns)
during the phase of adiabatic shortening prior to the fast rotation. This
again suggested that the short bunches could initiate the electron cloud build
up earlier in the cycle and produce enough electron cloud for a sufficiently
long time as to render the beam visibly unstable. In March 2007, an experiment
for dedicated electron cloud measurements was set up at the PS to be able to
directly measure the electron signal by using a shielded biased pick up [3]
and confirm its presence in the machine in the last phase of the LHC beams
production. The experimental setup was designed, fabricated, and mounted in
straight section (SS) 98 during the accelerator shutdown 2006/2007. All the
details of the setup can be found in Ref. [3]. These studies confirmed that
the electron cloud develops during the last 40 to 50 ms before ejection, i.e.
when the bunches are shortened by the RF gymnastics. Besides, they also showed
that the electron cloud can be suppressed by putting a sufficiently large
voltage of either polarity onto a clearing electrode, even if the clearing
efficiency depends on the magnetic field present in the region of the
measurement in a non-trivial way.
Table 1: Relevant beam parameters in the PS during the flat top RF gymnastics for the two bunch spacings of 50 and 25ns. | 50ns | 25ns
---|---|---
Beam energy (GeV/c) | 26
| Bunch intensity
---
($\times 10^{11}$ ppb)
0.82-1.95 | 0.83-1.33
Bunch length (ns) | 15 $\rightarrow$ 12 $\rightarrow$ 4
Number of bunches | 36 | 72
Transv. norm. emittances ($\mu$m) | 1-2 | 2-3
In 2011, new systematic measurements of electron cloud have been performed at
the CERN-PS with the goal of extracting the following information:
* •
Dependence of the electron cloud build up evolution on some controllable beam
parameters (bunch spacing, bunch intensity, bunch length).
* •
A new collection of time resolved experimental data of electron cloud build up
in some desired sets of beam conditions.
These sets of data can serve two purposes. First, comparing them with build up
simulations will allow us to validate (or improve) the simulation model on
which our tools are based. Second, by matching the simulations to the
experimental data in all the different beam conditions, we can pin down the
surface properties of the PS vacuum chamber (secondary electron yield,
$\delta_{max}$, and reflectivity of the electrons at zero energy, $R_{0}$) and
extrapolate then how much electron cloud we can expect in the PS with the
higher intensity beams foreseen in the frame of the LHC Injector Upgrade (LIU)
project, and whether that can be detrimental to the beam.
Table 2: Relevant beam parameters of the SPS 50 and 25ns beams. | 50ns | 25ns
---|---|---
Beam energy (GeV/c) | 26 $\rightarrow$ 450
| Bunch intensity
---
($\times 10^{11}$ ppb)
0.3-1.7 | 0.3-1.4
Bunch length (ns) | 4 $\rightarrow$ 2.8 $\rightarrow$ 1.6
Number of bunches | 144 | 288
Transv. norm. emittances ($\mu$m) | 1-2 | 2-3
The SPS has been suffering from electron cloud formation since it first began
to take and accelerate 25ns beams produced in the PS with the scheme explained
above. Observations of pressure rise, beam instability, emittance growth were
first made in the early 2000 and all these effects strongly limited the
capability of this accelerator of handling LHC-type beams [4]. While the
coherent instabilities could be suppressed by the use of the transverse damper
(against the horizontal coupled bunch oscillations) and running with
sufficiently high chromaticity (against the strong single bunch effect in the
vertical plane), emittance growth and positive tune shift along the bunch
train could still be measured, pointing to the continuing presence of a strong
electron cloud inside the beam chamber. All this led to the decision to have
in 2002 the first dedicated scrubbing run, in which the SPS was operated
exclusively with 25ns beams for one full week. The goal was to use the
bombardment from the electron cloud itself to clean the beam chamber inner
surface, and therefore lower its Secondary Electron Yield (SEY) and reduce, in
turn, the amount of electron cloud build up. The strategy proved successful
[5] and the week of scrubbing run was then repeated at the beginning of the
2003, 2004, 2006 and 2007 runs to provide the necessary machine cleaning.
During these years, dedicated experiments were conducted in the SPS to study
in detail the electron cloud formation in cold regions (COLDEX) or in NEG
coated chambers [6], or to benchmark simulation codes with machine
observations [7]. From 2006 on, electron cloud studies in the SPS acquired new
momentum in the framework of the SPS upgrade studies [8] and the experimental
activity over the following years was mainly focused to find the scaling law
of the electron cloud instability with beam energy [9] and to validate the
efficiency of amorphous carbon (a-C) coating of the beam chamber [10].
All the electron cloud machine development activity of the last couple of
years at the SPS has been devoted to defining the status of the 25ns beams in
this machine and use the direct electron cloud measurements in chambers
equipped with strip monitors to understand beam induced scrubbing in different
chamber geometries and with different materials. A comprehensive report of all
observations in terms of beam behavior, pressure rise and dedicated ekectron
cloud measurements will be soon published [11]. Recently, the nominal 25ns and
50ns LHC beams in the SPS seem to be undegraded and does not suffer any longer
from the strong electron cloud effect that was present during the first years
of SPS operation with this type of beams. The achievable parameters are
summarized in Table 2. The three values of bunch length quoted in this table
correspond to injection into 2 MV buckets, after shortening at flat bottom by
increase of the RF voltage to 3 MV, and at flat top after controlled
longitudinal emittance blow up during the accelerating ramp.
Several studies conducted in the past predicted that also the LHC would suffer
from heat load, pressure rise and beam instabilities due to electron cloud,
when operating with trains made of closely spaced proton bunches (e.g. [12]).
Since mid 2010 LHC entered this mode of operation. In the first phase, beams
with 150 ns bunch spacing were injected, accelerated and brought to collision.
During this period of operation, the only possible signature of electron cloud
build up was a pressure rise observed in the common vacuum chamber, close to
the Interaction Regions. Subsequently, at the end of October 2010, an attempt
was made to switch to 50 ns spacing operation. After an initial physics fill
with 108 nominal bunches (filling scheme with 1 pilot bunch and $9\times 12$
bunches), some important dynamic pressure rises were observed at injection
when filling with trains of 24 bunches. In fact, the first attempt of
injection in batches of 24 even led to the closure of the vacuum valves in
point 7 after the injection of 108 nominal bunches per beam, as the interlock
level of $4\times 10^{-7}$ mbar was reached on two vacuum gauges. After that,
since it became clear that further improvements in the LHC performance were
hampered by the electron cloud, emphasis was put on machine studies to
characterize the electron cloud build-up in the LHC, its effects and possible
cures. It was also decided that a comparative study with the behaviour of 75
ns beams was necessary to define a path for the 2011 run. Toward the end of
the 2010 proton run, a Machine Development (MD) session was devoted to the set
up of the LHC with 50 ns bunch trains. During this MD, three effective days of
beam time were used for the setting-up proper as well as for studies and
machine scrubbing. The study of the 75 ns beam took place in another dedicated
MD period, while the LHC had already switched to ion operation. About 2.5 days
were devoted to the setting-up of the injection and capture of the 75 ns beam
and, later on, to comparative studies with the 50 ns beam. This MD gave a
clear indication that, probably also benefiting from the previous MD’s
scrubbing with 50 ns beams, the electron cloud effects with 75 ns appeared
significantly less pronounced than with 50 ns beams, such that this bunch
spacing could be regarded as a relatively safe option [13].
Figure 1: MD sessions labeled (a), (c), (d) and (e): injected beams.
The LHC operation was therefore resumed in 2011 directly with 75 ns beams.
After the scrubbing run in 2010 it was expected that up to 200-300 bunches
could be injected and accelerated without major problems. This was confirmed
during the start-up with beam. After about one month of operation, the LHC
could successfully collide trains of 200 bunches distributed in batches of 24
bunches each. At the beginning of April, 10 days were devoted to scrubbing of
the LHC with 50 ns beams. The goal was to prepare the machine to switch to 50
ns beams and thus extend the luminosity reach for the 2011 run. During the
scrubbing run, up to 1020 bunches per beam were injected into the LHC in
batches of 36 and stored at injection energy. The strategy consisted of
constantly topping the total beam intensity in the LHC with the injection of
more trains, such that the vacuum activity, and therefore the electron cloud,
could be kept at a constant level and efficiently reduce the Secondary
Electron Yield (SEY) of the walls to a value below the threshold for build up.
The success of the scrubbing run was proved by the subsequent smooth LHC
physics operation with 50 ns spaced beams. Between mid April and end June the
number of bunches collided in the LHC was increased up to its maximum value of
1380 per beam, while the intensity per bunch and the transverse emittances
remained constant at their nominal values (i.e., 1.15$\times 10^{11}$ ppb and
2.5 $\mu$m). The switch to 50 ns beams with lower transverse emittances (1.5
$\mu$m) and the adiabatic increae of the bunch current to $1.5\times 10^{11}$
ppb did not cause any significant recrudescence of the electron cloud effects,
probably also owing to the MD sessions with 25ns beams that took place in the
second half of 2011, which created enough margin in the machine cnnditioning
to ensure electron cloud free 50ns operation.
Beams with 25ns spacing were injected into the LHC only during five MD
sessions of the 2011 run, which are listed and briefly described here:
1. (a)
29 June, 2011: first injections of 25ns beams into the LHC. The filling scheme
consisted of nine batches of 24 bunches separated by increasing gaps (2.28,
5.13 and 29.93 $\mu$s). Pressure rise around the machine as well as heat loads
in the arcs were observed. All the last bunches of each batch suffered losses
and emittance growth [14];
2. (b)
26 August, 2011: first injections of a 48-bunch train into the LHC with 25ns
spacing. Two attempts were made to inject a 48-bunch train from the SPS, which
led to beam dump triggered by large beam excursion and beam loss interlocks,
respectively. During the first injection test, the transverse damper was on
and it is believed that the beam suffered a coherent electron cloud
instability in both planes (more critical in vertical) soon after injection.
During the second test, the transverse damper was switched off and the beam
was affected by a coupled bunch instability [16]. This MD session had then to
be interrupted because of a cryo failure caused by a thunderstorm;
3. (c)
7 October, 2011: injection tests and first ramp. In the first part of the MD,
trains with 48-72-144-216-288 bunches from the SPS were injected into the LHC.
Given the experience during the previous MD, the chromaticity $Q^{\prime}$ was
set to around 15-20 units in both the horizontal and vertical planes in order
to keep the beams stable against the electron cloud effect. In the second
part, only 60 bunches per beam were injected in trains of $12+2\times 24$,
were accelerated to 3.5 TeV and collided during approximately 5h;
4. (d)
14 October, 2011: first long stores of 25ns beams at injection energy in the
LHC. During this session up to 1020 bunches per beam were injected in batches
of 72. The chromaticity was kept high in both planes ($Q^{\prime}_{x,y}\approx
15$) in order to preserve the beam stability. First, a dedicated fill for
pressure measurements was made, with batches injected at gradually reduced
distances from 4 to 2 $\mu$s (in steps of 1$\mu$s). Subsequently, the batch
spacing was kept constant for each of the next three fills and it was set to
6.3, 3.6 and 1 $\mu$s (rounded values). Strong emittance growth and slow
losses affecting the last bunches of each train were observed throughout this
MD session;
5. (e)
24–25 October, 2011: record number of bunches in the LHC. Four long fills took
place (average store time was approximately 4h), with 25ns beams injected into
both rings in batches of 72 separated by 1$\mu$s. In the third and fourth
fills, 2100 bunches were injected for beam 1, while the number of bunches
could not exceed 1020 for beam 2, due to a vacuum interlock on one of the
injection kickers (MKI). Although the situation seemed to improve over the MD,
slow losses and emittance growth kept affecting both beams. Before starting
the fourth fill, the horizontal chromaticity $Q^{\prime}_{x}$ was lowered from
15 to 3 units and the horizontal damper gain was slightly increased. Probably
due to that, some horizontal instabilities could be observed from the signal
of the damper pick up during the fourth fill, but the overall performance did
not appear degraded from the previous fill. The MD ended with a 30’ fill with
only beam 1, during which batches of 72 bunches were injected into the LHC at
different spacings in order to provide the stable pressure measurements needed
for the modeling of the electron cloud build up in the straight sections (see
next Section).
Figure 1 shows the detailed story, in terms of injected beams 1 & 2, of the
sessions (a), (c), (d) and (e). Experimental data from these MDs will be used
in the next section to extrapolate the evolution of $\delta_{\mathrm{max}}$ on
the beam screen in the arcs and in proximity of the vacuum gauges. For sake of
compactness, we have chosen to concatenate these three sessions and represent
them as a function of a continuous time coordinate (interpretable as hours
with 25ns beam), which will be systematically used throughout this paper when
referring to the studies with the 25ns beams.
## 2 Studies in the different machines
### 2.1 PS measurements
In 2011, the MD program in the PS for electron cloud studies took place in
November and extended over several sessions to cover different sets of beam
parameters. In particular, electron cloud build up data were recorded for 25ns
and 50ns beams. The bunch intensities were scanned in the ranges indicated in
Table 1. The trigger for the data acquisition was set at extraction, when in
normal conditions each bunch of the beam has been already fully rotated (4ns
bunch length). However, specifically for these measurements, the bunch length
at this time for a fixed bunch intensity was also set to 6.5ns or 15ns by
simply adjusting or fully removing, respectively, the final step of the fast
bunch rotation. This allowed studying the dependence of the electron cloud
build up not only on the bunch intensity but also on the bunch length.
The threshold for electron cloud formation with 50ns beams was found to lie at
about $10^{11}$ ppb and the measured signal increases monotonically with the
bunch intensity. This is not entirely surprising, since the measurements were
taken with zero magnetic field while the non-monotonic behaviour of the
electron cloud build up with the bunch intensity is more frequent in dipole
regions. The shielded pick up is installed inside a C-magnet, which was kept
off during the MD sessions because the orbit perturbation it introduces would
have required a specific correction. Scans with 25ns beams were also made and
the threshold for electron cloud formation was found to be below $8\times
10^{10}$ ppb, with a behavior of the electron cloud signal increasing with the
bunch intensity.
We have tried to fit the PS data with those from electron cloud build up
simulations [17]. First of all, the output of the code that should be compared
with the measured signal is the electron flux to the wall. In a first
approximation, we do not consider the holes in the vacuum chamber, which are
expected to cause only a minor perturbation in a field-free region. In
general, the simulated electron flux to the wall vanishes during the bunch
passage, because initially all the electrons are drawn to the center of the
vacuum chamber by the passing bunch (e.g. during the first $\sim$2ns of a 4ns
long bunch) and they are gradually released only during the falling edge of
the bunch, when they may reach the walls again. The fact that the measured
signal does not exhibit this feature makes plausible a low pass filtering of
the signal (inherent to the measurement technique or due to electronics and/or
cables) with a corner frequency in the range of some hundreds of MHz. Figure 2
shows measured and simulated signal, where the simulated signal, obtained with
$\delta_{\mathrm{max}}=1.6$ and $R_{0}=0.5$, was low pass filtered with a
corner frequency of 200 MHz. The impressive resemblance between the two
suggests that our electron cloud model correctly describes the phenomenon and
the rationale applied for the data analysis is promising.
Figure 2: E-cloud build up simulation (top) and measurement (bottom) for a
25ns beam with $1.33\times 10^{11}$ ppb and 4ns long.
It is therefore clear that the electron cloud is present in the CERN PS with
both 50ns and 25ns beams when they reach the final beam structure, shorty
before being ejected. However, since it only makes a short appearance in the
last few ms of the production cycle of these beams, with the present beam
parameters, there is not enough time to render beam unstable or let incoherent
effects develop. On the negative side, very low electron doses are deposited
on the chamber walls, making it basically impossible to rely on efficient
machine scrubbing if the electron cloud should ever become a limiting factor.
The question to be addressed is whether this effect may become a bottleneck
for the LHC Injector Upgrade (LIU) beams, envisaging bunch intensities of up
to $3\times 10^{11}$ ppb within lower transverse emittances. A full simulation
study including both the build up and instability part is needed to assess the
margins.
Figure 3: Threshold SEY for electron cloud formation in the two types of
chambers in SPS drift spaces, as a function of the bunch intensity.
### 2.2 SPS studies
One of the key points to be addressed to understand the electron cloud in the
SPS is to determine the values of SEY thresholds for its formation in the
different beam chambers and try to deduce what parts are critical for both
present and future LHC beams. In the SPS there are six main different types of
vacuum chamber: two types are used in the main dipoles, two in the quadrupoles
and two in the drift spaces, depending on the beta functions in the nominal
optics. We have studied the electron cloud build up in both dipole chambers
and driift space chambers. Since the magnetic fields from quadrupoles have not
been implemented yet in the PyECLOUD code, these chambers, which however cover
only less than 10% of the total circumference, have not yet been simulated.
The drift chambers are of A or B type, both circular and with a radius of 78
or 65 mm, respectively. Figure 3 shows the SEY threshold as function of bunch
intensity at both 26 and 450 GeV/c. The following interesting features can be
observed:
* •
The SEY thresholds are mostly decreasing with bunch current, but tend to
change slope for 50ns beams with bunch populations above $2\times 10^{11}$
ppb.
* •
There exist regions in which 50ns can create a worse electron cloud than 25ns.
* •
The SEY thresholds become very low (close to 1.05) for 25ns beams in Drift B
and with bunch currents above $2\times 10^{11}$ ppb.
The vacuum chambers in dipoles also come in two different sorts with almost
rectangular shape: the MBB-type, characterized by a height of about 56.5mm and
132mm width; and the MBA-type, flatter than the MBBs and thus more suited to
regions with lower vertical betatron functions, characterized by a height of
43mm and 156mm width. In Fig. 4 the SEY threshold is displayed as a function
of the bunch intensity for both 25 and 50ns beams as well as at injection and
top energy. Also in this case the dependencies are not trivial and exhibit the
following features:
* •
The SEY thresholds do mostly increase with bunch current. When they do not,
the behavior tends to be flat, indicsting then little dependence of the SEY
threshold on the bunch intensity in these intensity ranges.
* •
The SEY thresholds of the 50ns beam lie above 2.0 in the MBA chambers.
* •
The SEY thresholds can become in general very low (around 1.2) for 25ns beams
in MBB chambers.
Figure 4: Threshold SEY for electron cloud formation in the two types of
chambers in SPS dipoles, as a function of the bunch intensity.
Considering all the results of the above study, it is evident that the most
dangerous chambers in the SPS, in terms of favoring electron cloud build up
for the present and future LHC beam intensities, are the drift B and the MBB
pipes, which exhibit the lowest SEY thresholds at almost all intensity ranges.
In particular, it is specially worrisome that both these chambers exhibit SEY
thresholds below 1.3, which seems the saturation value for scrubbing of StSt
in laboratory measurements [18]. Besides, StSt samples exposed to the SPS beam
and then extracted from the machine have never shown SEY values below 1.5.
Presently, it still remains unclear whether we still have electron cloud in
some of the SPS regions, because the observed pressure rise is several order
of magnitude lower than the one observed in previous years and the nominal
25ns beam is not really affected anymore by significant electron cloud effects
[11]. While more studies are ongoing to try to characterize the present status
of the SPS and draw conclusions on future strategies against electron cloud,
it is however clear that the critical regions that might need coating (if
scrubbing is insufficient or too long) would amount to about 40% of the whole
machine (Drift B + MBB).
Experimentally, we can say that, thanks to the regular scrubbing runs the SPS
had from 2003 to 2008 with 25ns beams at every start up (plus several MD
sessions with this type of beams every year), the performance with 25ns beams
has been constantly improving over the years and in 2011, nominal 25ns beams
with transverse emittances below 3$\mu$m were first produced and extracted.
This leads us to believe that presently the electron cloud has weakened or
disappeared in most parts of the SPS and might be still only surviving in the
MBBs for operation with nominal intensity 25ns beams. In these conditions, it
seems to be efficiently kept under control and does not give rise to
detrimental effects on the beam. An Increase in bunch intensity may awaken the
electron cloud in the Drifts (and MBAs, because the stripes move to unscrubbed
regions) with the consequent effect on beam stability and emittance evolution.
It is clear that a scrubbing run will be necessary after the Long Shutdown
2013-2014, but its length and efficiency are difficult to estimate. The
experience after LS1 will therefore give an indication whether we really need
to coat the most critical parts of the SPS, or we can afford to rely on
scrubbing also for future operation.
Figure 5: Top picture: Heat load measured during four fills from the MD
session (d) and (e), in the same time coordinate as in Fig.1. The five
vertical bars represent the measurement points used to compare heat load with
electron cloud simulations. Bottom picture: Estimated evolution of
$\delta_{\mathrm{max}}$ on the inner surface of the beam screen in the dipole
chambers
### 2.3 LHC observations
The heat load data from the cryogenic system give the total power dissipated
(in W/half-cell) on the beam screens of both beams 1 and 2. Using the measured
heat load it is possible to estimate the SEY of the arc chamber walls. The
exact procedure is explained in Ref. [15] and is based on the comparison of
the heat load data with PyECLOUD simulations, run with realistic bunch-by-
bunch intensities and lengths (data from the fast BCT and the BQM). Heat load
observations in the arcs were made with 50ns before the scrubbing tun and then
with 25ns beams. From the heat load data with 50ns beams before and after the
scrubbing run, we could estimate the SEY on the chamber wall of the arcs to
have reached a value between 2.1 and 2.2, sufficient to suppress electron
cloud build up with 50ns beams. Measurements in some reference cells from the
first LHC MD with 25ns beams (MD session (a), 29 June, 2011) can be found in
Ref. [14]. Figure 5 shows the heat load data, sector by sector, collected
during the MD sessions (d) and (e). We can notice that the additional heat
load peaked to values of nearly 50 W/half-cell (i.e. approximately an average
of 0.5 W/m/beam) during the last fill with 2100 bunches for beam 1 and 1020
bunches for beam 2. A decay of the measured heat load between injections, and
in any case after the last injection, is also clearly visible in the examined
cases, due to the weakening of the electron cloud activity from scrubbing and
also from intensity loss (e.g., compare with the BCT signal in Fig. 1,
acquired at the same time).
Using the bunch-by-bunch intensity and length data at the times marked with
vertical bars in the top plot of Fig. 5 plus the data from the injection tests
on the 29 June, PyECLOUD simulations were run scanning
$\delta_{\mathrm{max}}$, so that the curves of the simulated heat loads as a
function of $\delta_{\mathrm{max}}$ could be produced for all these
measurement points. The electron reflectivity at zero energy was fixed to the
value of 0.7. The $\delta_{\mathrm{max}}$ corresponding to each heat load
measurement was then found matching the simulation to the measured value and
the results are in the curve displayed in the bottom part of Fig. 5.
Figure 6: Snapshots of the horizontal and vertical emittance measurements for
beam 1 during the last fill of 14 October and the last two fills of 24–25
October MDs.
While the 50ns beam proved to be stabilized in the LHC by the electron cloud
mitigation achieved with the scrubbing run, the 25ns beam has exhibited clear
signs of transverse instability and emittance growth throughout all the
dedicated MD sessions. Despite a clearly improving trend from one fill to the
next one, these signs have not completely disappeared. During the first tests
on 29 June, when only batches of 24 bunches were injected from the SPS, the
beam could be kept inside the machine because the level of electron cloud
reached along each batch was enough to cause significant emittance growth, but
no coherent instability and fast beam loss [14]. When, on the following MD
session, batches of 48 bunches were for the first time transferred from the
SPS to the LHC, the beam was twice dumped after few hundreds of turns, due to
the excitation of a transverse instability leading to unacceptable beam
losses. During the successive MD sessions, this problem was circumvented by
injecting the beam into the LHC with high chromaticity settings. Values of
$Q^{\prime}_{x,y}$ around 15 were chosen, as they had been found to be
sufficiently stabilizing in HEADTAIL simulations [19]. Using these settings,
the beam could be kept inside the LHC, albeit with degraded transverse
emittances (see bunch-by-bunch emittance plots from the MDs of 14, 24 and 25
October, Fig. 6). Since the BSRT needs about 2 sec to measure the emittances
of each bunch, each of the snapshots in the figure does not represent an
instantaneous photograph of the beam at a certain time, but results from a
sweep over the bunches that can last as much as several minutes. Although the
batch spacing was decreased from 2 $\mu$s during the measurement of 14 October
to the 1 $\mu$s of the last MD session, the vertical emittance blow up
exhibits signs of improvement. No significant further change is observed then
in the vertical plane between the measurements taken in the last two fills
(consistently with a slight scrubbing effect between them). The situation
looks more complicated in the horizontal plane. Here a deterioration can be
noticed from the 14/10 measurement to the 24/10 one. If this is related solely
to the decreased batch spacing, which has enhanced the electron cloud along
the full train owing to the stronger memory effect between batches, we could
not explain why we observed an improvement in the vertical plane, instead. It
is interesting that the situation appears improved for the 25/10 measurement,
when the LHC was run with lowered horizontal chromaticity settings. This fact
may suggest that by lowering chromaticity we have moved from a regime of
strong incoherent emittance growth driven by electron cloud and high
chromaticity to a new one, in which the beam suffers a fast instability, but
later evolves with a better lifetime [20, 21]. In any case, as a general
consideration, a clear weakening of the electron cloud effect from 14 to 25
October is witnessed by the improved quality of the first two–three batches.
The first two seem to be hardly affected by emittance growth in both
transverse planes by the time of the last 25ns fill.
## 3 Conclusions and outlook
In conclusion, we have reached quite a deep knowledge of the electron cloud in
the different CERN accelerators, and presently it does not seem to be a
limiting factor with the present operation parameters:
* •
In the PS the electron cloud only appears in the last milliseconds of the LHC
beam production cycle and does not stay long enough as to affect the beam
* •
The SPS currently benefits from several years of scrubbing with nominal 25ns
beams. Therefore, it seems that now the electron cloud has been either
suppressed in the whole machine or it still survives in some more sensitive
parts, but at a level not harmful to the beams (50ns, nominal 25ns)
* •
In the LHC, the electron cloud does not have important adverse effects on
operational 50ns beams, however it still affects the 25ns beams and additional
scrubbing is needed to further lower it and permit operation with this type of
beams.
However, some questions are still open, like the performance of SPS and LHC
with 25ns beams after LS 2013-2014 and whether the electron cloud can become a
serious bottleneck for the beams required by the LIU project.
## References
* [1] R. Garoby, “Status of the Nominal Proton Beam for LHC in the PS”, CERN/PS 99-13 (RF)
* [2] R. Cappi, et al., Phys. Rev. ST Accel. Beams 5, 094401 (2002).
* [3] E. Mahner, T. Kroyer and F. Caspers, Phys. Rev. ST Accel. Beams 11, 094401 (2008)
* [4] G. Arduini, K. Cornelis, W. Hoefle, G. Rumolo, and F. Zimmermann, in Proceedings of PAC 2001 (18-23 June 2001, Chicago, USA) and CERN-SL-2001-0050
* [5] J.M. Jimenez et al., LHC-Project-Report-632 (2003)
* [6] V. Baglin, A. Rossi, et al. in Proceedings of ECLOUD04 (19-23 April 2004, Napa California, USA)
* [7] D. Schulte, G. Arduini, V. Baglin, J.M. Jimenez, F. Ruggiero, and F. Zimmermann , in Proceedings of PAC2005 (16-20 May 2005, Knoxville Tennessee, USA) and LHC Project Report 847
* [8] G. Rumolo, E. Métral and E. Shaposhnikova, in Proceedings of LHC LUMI 2006 (16-20 October, Valencia, Spain)
* [9] G. Rumolo, G. Arduini, E. Métral, E. Shaposhnikova, E. Benedetto, R. Calaga, G. Papotti and B. Salvant, Phys. Rev. Letters 100 (2008) 144801
* [10] C. Yin Vallgren et al., Phys. Rev. ST Accel. Beams 14, 071001
* [11] H. Bartosik, G. Iadarola et al., to be published
* [12] F. Zimmermann, in Proceedings of Chamonix X & XI, CERN-SL-2000-001 DI (2000) and. CERN-SL-2001-003 DI (2001)
* [13] G. Arduini et al., CERN-ATS-Note-2011-046 MD (2011)
* [14] B. Goddard et al., CERN-ATS-Note-2011-050 MD (2011)
* [15] G. Rumolo et al., in Proceedings of the LHC Beam Operation Workshop - Evian 2011 (12-14 December, 2011, Evian, France)
* [16] H. Bartosik and W. Höfle, CERN-ATS-Note-2012-027 MD (2012)
* [17] G. Iadarola and G. Rumolo, “Improved electron cloud build up modeling with PyECLOUD”, elsewhere in these proceedings
* [18] C. Yin Vallgren, Ph.D. thesis, CERN-THESIS-2011-063 (2011)
* [19] K. Li and G. Rumolo, MOPS069 in proceedings of IPAC’11 (San Sebastian, Spain)
* [20] E. Benedetto, G. Franchetti, and F. Zimmermann, Phys. Rev. Lett. 97 (2006) 034801
* [21] K. Ohmi and F. Zimmermann, Phys. Rev. Lett. 85 (2000) 3821-3824
|
arxiv-papers
| 2013-09-26T11:03:55 |
2024-09-04T02:49:51.561276
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "G. Iadarola (CERN and U. Naples), G. Rumolo (CERN)",
"submitter": "Scientific Information Service CERN",
"url": "https://arxiv.org/abs/1309.6795"
}
|
1309.6914
|
# Chemical concrete machine
Marius Buliga
Institute of Mathematics, Romanian Academy
P.O. BOX 1-764, RO 014700
Bucureşti, Romania
[email protected]
(This version: 26.09.2013)
###### Abstract
The chemical concrete machine is a graph rewriting system which uses only
local moves (rewrites), seen as chemical reactions involving molecules which
are graphs made up by 4 trivalent nodes. It is Turing complete, therefore it
might be used as a model of computation in algorithmic chemistry.
## 1 Introduction
Moe-Behrens explains [10] that there is a trend to use the architecture of a
silicon computer for building a biological computer. Whatever happens in a
living cell which resembles to computation (in a more vague sense than Turing
computation) is not necessarily implemented as in a silicon computer, though.
Instead, synthetic biology tries to import concepts from silicon computers
into the bio realm, mainly because these concepts are already familiar. Hence
the accent on bits and boolean logic, although not exclusive (other studies
concern for example membrane computing, or neural networks, or continuously
evolving dynamical systems). More important, all these studies seem to
concentrate on imperative programming.
Functional programming is the other main programming paradigm, a very elegant
one, receiving more and more attention in the competition with the more well
known imperative programming. Functional programming has its roots in the
lambda calculus, one of the two pillars of computation, along with the Turing
machine (of equivalent power, but with different philosophies behind).
In 1992 Berry and Boudol [2] introduce the chemical abstract machine, which
models asynchronous concurrent computations (in particular a concurrent lambda
calculus) by using a chemical metaphor. The chemical abstract machine has
states which are seen as chemical solutions and the evolution from a state to
another is ruled by chemical reactions, heating and cooling. Chemical
solutions are multisets of molecules, i.e. list of numbers of molecules of
each type.
Between 1994-1996, Fontana and Buss [6] [7] [8] introduce the idea that lambda
calculus is a kind of natural formalization of the bare bones of chemistry.
Individual molecules are seen as lambda terms, reactions between molecules are
seen as the application operation in lambda calculus, and the abstraction
operation from lambda calculus ”captures the role of functional groups
embedded in the reactants”, [7] p. 11. Fontana and Buss use lambda calculus
with eta reduction, i.e. the terms (or molecules) are identified with
functions, in the sense that the molecule $A$ is identified with the function
$B\mapsto AB$.
The work of Fontana and Buss started the ”algorithmic chemistry” research
field. In a sense, this can be seen as a part of an older research subject,
the very difficult one, of course, concerning the organization and functioning
of living systems. They cite in [7] the seminal paper by Varela, Maturana and
Uribe [12].
The chemical concrete machine is a graphic formalism which uses a chemical
metaphor. It is a modification of graphic lambda calculus, introduced in [3].
(See also the graphic lambda calculus web tutorial. Moreover, there is already
a web tutorial for the chemical concrete machine.)
The graphic lambda calculus has been constructed for reasons independent from
any of the subjects evoked until now. It is a graph rewriting system which has
a sector which is equivalent with untyped lambda calculus, but it has other
interesting sectors as well, like the one which is equivalent with knot
diagrams, or the one which contains finite differential calculus in spaces
with dilations [4], i.e. emergent algebras [5]. It is therefore Turing
complete, but in order to achieve this, it uses a GLOBAL FAN-OUT move. Moves
(i.e. graph rewrites) can be characterized as local, when they involve a fixed
number of nodes and arrows in the graphs, or global, otherwise.
By modifying the graphic lambda calculus, we obtain the chemical concrete
machine, which is Turing complete by using only local moves. Another
difference from graphic lambda calculus is that we reformulate the moves as
chemical reactions between molecules (represented by trivalent graphs made by
4 types of trivalent nodes). Thus a computation in the chemical concrete
machine is simply a chemical reaction network involving molecules and chemical
reactions from the formalism.
There are also some differences between the algorithmic chemistry of Fontana
and Buss and the chemical concrete machine. The most important one is that
that the application and the abstraction operations from lambda calculus are,
in the chemical concrete machine, trivalent nodes (say, like atoms which
compose the molecules), and moves (corresponding to reductions in lambda
calculus) are assimilated with enzymes, thus reductions are certain chemical
reactions. Also, the chemical concrete machine can be used for implementing
lambda calculus without extensionality, mainly because the eta reduction
appears as a global move in graphic lambda calculus, which does not seem
possible to be replaced by local moves, as it has been done with the GLOBAL
FAN-OUT.
This leads to the last point: it is intriguing to search if ”extreme”
functional programming (i.e. functional programming without extensionality) is
easier to be implemented in biological computing than imperative programming.
Already, lambda calculus looks like a kind of chemistry. In the chemical
concrete machine the somehow esoteric application and abstraction are made
concrete. Moreover, because the chemical concrete machine uses the graphic
lambda calculus formalism, it follows that simple spatial manipulations, like
grouping molecules together, releasing them in some order or all of them
simultaneously, are more straightforward to achieve in graphic lambda calculus
than in the usual lambda calculus.
Another possibility would be to try to implement the chemical concrete machine
in silico, which will be discussed in another article.
### Acknowledgments.
I benefited from various discussions with the following (in no particular
order): Stephen P. King, Eugenio Battaglia, Gerd Moe-Behrens, Louis Kauffman,
ine (nickname in comments made at the chorasimilarity blog), Lucius Meredith,
Mike Stay.
## 2 Chemical concrete machine
The chemical concrete machine is a graph rewriting system, derived from
graphic lambda calculus [3], which can be seen as a model of chemical, or why
not biological, computing.
The distinction from other such models is that it works with very specific
graphs, called ”molecules”, which interact in very specific ways. This is the
source of the word ”concrete” in the denomination, as opposed to the chemical
abstract machine of Berry and Boudol [2].
I hope it can be implemented in reality (that is in a real chemical reaction
network), but nevertheless it can be seen as a proof of principle for the fact
that the chemistry of a very small number of molecules can manifest Turing
completeness. In this respect, the chemical concrete machine appears as
applied algorithmic chemistry in a virtual world with chemistry rules which
are made up by a mathematician.
Molecules. In the following I shall use the name ”molecule” in the following
sense. A molecule is seen as made by other smaller molecules (for example
atoms are molecules) which are connected by chemical bonds, or by other
molecules called ”arrows”. Therefore, a molecule is seen as a graph with nodes
which are smaller molecules and arrows which connect those nodes. A collection
of molecules is a molecule.
The building block of molecules are the following elementary, or essential
ones.
On the first row you see the list of essential molecules, named respectively
1. -
”application”, corresponding to the application gate in graphic lambda
calculus, or to the application operation in lambda calculus
2. -
”fan-out” , corresponding to the FAN-OUT gate in graphic lambda calculus
3. -
”fan-in”, which will replace the $\varepsilon$ gate from graphic lambda
calculus
4. -
”abstraction”, corresponding to the $\lambda$ gate in graphic lambda calculus,
or to the lambda abstraction from lambda calculus.
On the second row we see a loop, an arrow (which are considered molecules, as
written before), and a molecule called ”terminal”, which corresponds in
graphic lambda calculus to the termination gate.
We shall work also with a list of ”other molecules”, which can react (or not)
one with another, this is left to the choice of the user of the chemical
concrete machine, which can add more reactions in the formalism. There is also
a collection $GARB$ of garbage.
###### Definition 2.1
The set $MOLECULES$ of molecules is formed by all locally planar graphs which
can be formed by using the essential molecules (and the ”other molecules” with
unspecified valences) as nodes, and arrows obtained by connecting the
essential molecules such that they respect the arrows orientation. The graphs
can have free arrows (with one of the ends of an arrow not connected to any
other molecule), also they can have several connected components (i.e. a
finite union of molecules is a molecule). Loops and arrows are molecules.
The rules of building molecules are exactly the ones for building graphs in
the set $GRAPH$ described in section 2 [3], or go directly to the graphic
lambda calculus tutorial, with the understanding that for the chemical
concrete machine we use colored nodes instead of the gates from graphic lambda
calculus, as described in the following figure. The only new thing is that we
admit also a list $A,B,C,...$ of unspecified nodes with unspecified valences,
which model ”other molecules”.
Here are some examples of molecules.
We imagine that all molecules float inside a container. There may be several
copies of the same molecule in the container.
Enzymes and moves. Besides molecules, I shall also use ”enzymes”. A chemical
reaction will always involve a molecule and an enzyme. Enzymes are not
molecules. Instead, enzymes are names of moves (graph rewrites).
What is to be noticed is that any enzyme comes in two flavors: ”+” and ”-”.
The reason for this is the following. Any chemical reaction which involves a
molecule and an enzyme will correspond to a move in the list of moves of the
chemical concrete machine formalism. The list of moves of the chemical
concrete machine, which will be described soon, contains the local moves of
graphic lambda calculus which involves the fan-out, application and
abstraction gates, supplemented with a list of moves related to the fan-in
gate, which replace the emergent algebra moves from the graphic lambda
calculus. All the moves of the graphic lambda calculus are bi-directional.
Therefore, any such move can be seen as a pair of unidirectional moves. By
convention, the moves from left to right are denoted by names of enzymes with
”+” and the moves from right to left are denoted by names of enzymes with ”-”.
We cannot simply write chemical reactions as
molecules + enzyme = molecules + GARB
because these molecules may be complicated beasts (graphs) and because, as we
shall see, the enzymes prefer to react with certain patterns (subgraphs) of
molecules. For specifying how the reaction takes place we need:
1. -
molecules
2. -
an enzyme
3. -
and a ”reaction site”, which is a small part of the initial collection of
molecules.
The reaction site have to be present in the molecule, otherwise the reaction
cannot happen.
In order to explain how to translate from the language of chemical reactions
and enzymes to the language of moves of the chemical concrete machine, let’s
contemplate the following figure. I consider two examples of reactions (which
will turn out to correspond to graphic beta moves in the realm of graphic
lambda calculus).
In the first row we see a reaction between a molecule and the enzyme
$\beta^{+}$, which results into two other molecules and some GARB. There is a
small region in the initial molecule, marked by a dashed red closed curve,
which represents a reaction site for the $\beta^{+}$ molecule.
In the second row is written the same reaction, but in a simpler form. The red
”+” sign is eliminated, the two molecules which are obtained are juxtaposes,
as if they are floating in the 3D container, and the GARB is ignored.
Moreover, the enzyme $\beta^{+}$ points towards the reaction site.
The rows 3 and 4 describe another reaction. At closer inspection, it’s a
reaction which can be interpreted as the inverse of the first one. Let’s
examine directly the 4th row (which is obtained from the 3rd row by the same
procedure as the 2nd row was obtained from the 1st). The reaction site of the
enzyme $\beta^{-}$ is a pair of arrows from two different molecules, The
resulting molecule is the same as the initial molecule from the previous
reaction.
###### Definition 2.2
Here is the list of moves, with the names taken from graphic lambda calculus
(where it is the case), by using the dictionary for translation (links to the
web tutorials are given).
1. -
graphic beta move and FAN-IN move :
1. -
local FAN-OUT moves :
1. -
DIST moves (”dist” comes from the word ”distributivity”):
1. -
LOCAL PRUNING moves:
1. -
Elimination of loops (i.e. moving loops from and to $GARB$)
While all of these moves define the chemical concrete machine, some of them,
seen in isolation, appear in other parts. I don’t know if there is any, more
profound, reason for this, but nevertheless here are some examples of such
places:
1. -
in [1], the merge move (2.28) section 2.6 looks like the pair of graphic beta
move - fan-in move; the blow up move (13.1) section 13.1 looks even more like
the graphic beta move,
2. -
in [11], the unzip move from section 3 looks almost alike the graphic beta
move,
3. -
in [9], after definition 3.1.3 of a bialgebra in a monoidal category, a DIST
move and a LOC PRUNING move appear under the form of graphical identities. Of
course, CO-COMM and CO-ASSOC moves can be seen as related to co-commutative
and co-associative comonoids.
## 3 Using the chemical concrete machine (I)
Let’s see what the chemical concrete machine can do.
### Lists and locks.
Suppose you have a family of molecules which you want to free in the medium in
a given order. This corresponds to having a list of molecules, which is ”read”
sequentially. I shall model this with the help of the zipper from graphic
lambda calculus.
Suppose that the molecules we want to manipulate have the form $A\rightarrow
A^{\prime}$, with $A$ and $A^{\prime}$ from the family of ”other molecules”
and $\rightarrow$ an arrow. Here are three zippers.
The first zipper, called a $\beta$ zipper, behaves in the following way. In
the presence of $\beta^{+}$ enzymes, there is only one reaction site
available, namely the one involving the red and green nodes in the
neighbourhood of the $D,D^{\prime}$. So there is only one reaction possible
with a $\beta^{+}$ enzyme, which has a a result the molecule $D\rightarrow
D^{\prime}$ and a new, shorter $\beta$ zipper. This new zipper has only one
reaction site, this time involving nodes in the neighbourhood of
$C,C^{\prime}$, so the reaction with the enzyme $\beta^{+}$ gives
$C\rightarrow C^{\prime}$ and a new, shorter zipper.
The reaction continues like this, freeing in order the molecules $B\rightarrow
B^{\prime}$, then $A\rightarrow A^{\prime}$ and $E\rightarrow E^{\prime}$.
The second zipper is called a FAN-IN zipper (or a $\phi$ zipper). It behaves
the same as the previous one, but this time in the presence of the FAN-IN
enzyme $\phi^{+}$.
In the case of a mixed zipper, he first molecule $D\rightarrow D^{\prime}$ is
released only in the presence of a $\phi^{+}$ enzyme, then we are left with a
$\beta$ zipper.
This can be used to lock zippers. Look for example at the following molecule:
called a locked $\beta$ zipper. In the presence of only $\beta^{+}$ enzymes,
nothing happens, because there is no reaction site compatible with the
$\beta^{+}$ enzyme. . If we add into the reactor also $\phi^{+}$ enzymes, then
the zipper unlocks, by releasing a loop (that’s seen as garbage) and a $\beta$
zipper which starts to react with $\beta^{+}$ enzymes.
The same idea can be used for keeping a molecule inactive unless both
$\phi^{+}$ and $\beta^{+}$ enzymes are present in the reactor. Say that w have
a molecule $A\rightarrow A^{\prime}$ which is made inactive under the form
presented in the following figure
The molecule is locked, but it has two reaction sites, one sensible to
$\beta^{+}$, the other sensible to $\phi^{+}$. Both enzymes are needed for
unlocking the molecule, but there is no preferred order of reaction with the
enzymes (in particular these reactions can happen in parallel).
### Sets.
Suppose now that we don’t want to release the molecules in a given order. We
need to prepare a molecule which has several reaction sites available, so that
multiple reactions can happen in parallel, as in the last example.
Mathematically, that could be seen as a representation of the set of molecules
we want to free, instead of the list of them. This is easy, as described in
the next figure:
On the first row we see what is called a $\beta$ set. It has 4 possible
reaction sites with the enzyme $\beta^{+}$, therefore, in the presence of this
enzyme, the molecules $A\rightarrow A^{\prime}$, … , $E\rightarrow E^{\prime}$
are released at the same moment.
A FAN-IN, or $\phi$ set, behaves the same as the previous one, but this time
in the presence of the FAN-IN $\phi^{+}$ enzyme.
Finally a mixed set releases different molecules, depending on the presence of
$\phi^{+}$ or $\beta^{+}$ enzymes.
### Pairs.
As another example, here is a more involved molecule, which produces different
pairs of molecules, according to the presence of $\phi^{+}$ or $\beta^{+}$
enzymes.
In the following figure we see how we model a pair of molecules, then two
possible reactions a represented.
The idea is that we can decide, by controlling the amount of $\beta^{+}$ or
$\phi^{+}$, to couple $A$ with $D$ and $C$ with $D$, or to couple $A$ with $B$
and $C$ with $D$.
4\. Multipliers and comultipliers. Multipliers and co-multipliers are
molecules which self-multiply. More precisely, in the next figure we see the
definition of those:
Here $A$ and $A^{\prime}$ are molecules from the formalism of the chemical
concrete machine and $1$ and $2$ are labels. The blue arrow means any finite
sequence of chemical reactions (moves) from the formalism.
In certain configurations, zippers are multipliers. In the following figure we
see what happens in the presence of DIST enzymes:
The reaction continues:
Now, the zipper multiplied into two zippers, but they are still connected. We
need more information about $A,B,C,D$ and
$A^{\prime},B^{\prime},C^{\prime},D^{\prime}$. Remark that:
In conclusion: if $A,B,C,D$ are multipliers and
$A^{\prime},B^{\prime},C^{\prime},D^{\prime}$ are co-multipliers, then the
zipper is a multiplier. This will be used in the next section.
## 4 Using the chemical concrete machine (II). Turing universality
The B,C,K,W system ”is a variant of combinatory logic that takes as primitive
the combinators $B,C,K$, and $W$. This system was discovered by Haskell Curry
in his doctoral thesis”.
I shall explain first which are the correspondents of the B, C, K, W,
combinators in the formalism of the chemical concrete machine. (Via the red-
green vs black-white change of notation, they can be deduced from their
expressions in graphic lambda calculus, which are obtained by using the
algorithm described in section 3 [3]. )
###### Definition 4.1
The class of combinator molecules (i.e. those which correspond to combinators
written in the B,C,K,W system) is defined by the following procedure:
1. -
the B, C, K, W molecules are combinator molecules,
2. -
if A, D are combinator molecules then the molecule AD, defined as in the next
figure, is a combinator molecule,
3. -
if A is a combinator molecule and B is another molecule obtained from A by the
application of the moves of the chemical concrete machine, then B is a
combinator molecule.
The purpose of this section is to prove the following.
###### Theorem 4.2
The chemical concrete machine contains combinatory logic under the form of the
B,C,K,W system.
### Proof.
By Theorem 3.1 and Proposition 3.2 [3], using also the dictionary between the
notations in graphic lambda calculus and those of the chemical concrete
machine, it follows that the combinatory algebra (in particular in the form of
the B,C,K,W system) can be expressed in the chemical concrete machine if:
1. -
we replace the GLOBAL PRUNING move from graphic lambda calculus with the
management of the GARB (garbage) class of molecules (alternatively, we might
simply neglect GLOBAL PRUNING and we obtain something a bit more expressive
than combinatory algebra),
2. -
we add to the chemical machine list of moves the following GLOBAL FAN-OUT
move:
Here, the meaning of the red dashed circle is that the GLOBAL FAN-OUT move can
be applied only if the molecule A has no arrows outside the red dashed circle,
with the exception of the one connecting it with the fan-out molecule.
In Proposition 3.2 [3] was used the S,K,I combinators system, instead of the
B, C, K, W, but the Theorem 3.1 [3] shows that already untyped lambda calculus
can be expressed with the chemical concrete machine, provided that the GLOBAL
FAN-OUT move is added. The choice of the B,C,K, W system instead of the S,K,I
system is only for exposition purposes.
The goal is to show that for any combinator molecule, the GLOBAL FAN-OUT move
can be replaces by a finite chain of local moves of the chemical concrete
machine. As a consequence, we don’t need the GLOBAL FAN-OUT move in order to
have combinatory logic in the chemical concrete machine.
By looking at the last part of the previous section, we want to prove that any
combinator molecule is a multiplier. This is done in the following two steps.
Step 1. If A, D are combinator molecules which are multipliers then AD is a
multiplier. The proof is given in the next figure, by using DIST moves:
Step 2. Is only left to prove that B,C, K, W are multipliers. By the
conventions of the chemical concrete machine, I mention here the enzymes which
are involved in the reactions, instead of writing the moves, like in the
graphic lambda calculus.
The proofs for B and C are very much alike, therefore I put here only the
proof that B is a multiplier:
The proof that K is a multiplier is the following:
As a side remark, notice how, in both cases, the reactions seem feasible, in
the sense that they can be accomplished by a linear process, because at any
step there is only one kind of reaction site available.
For the W combinator (molecule), things get a bit more complex.
There is a reaction (or move) which needs explanations. I called it
DISENTANGLE (CO-ASSOC) reaction. It is this:
It can clearly be done by a succession of CO-ASSOC moves (reactions). This
ends the proof of the theorem. $\square$
From the point of view of the feasibility in the real world (provided a real
implementation of the chemical concrete machine will appear), it seems hard to
control the exact order of applications of CO-ASSOC moves which gives the
DISENTANGLE move as an effect. So, probably, we shall need a ”disentangle
enzyme” dedicated to this.
As an alternative, remark that for proving that W is a multiplier we need an
application of the DISENTANGLE composite move, described in the next figure:
For practical (or theoretical as well) purposes, it is enough to take this as
a move. In other words, what would get rid of needing a controlled sequence of
CO-ASSOC reactions for multiplying the molecule W is this: assume that the
molecule which is connected to the ”y” essential molecule (i.e. to the input
of a FAN-OUT gate) is a ”propagator”. Propagators are defined in the next
figure:
Propagators are different from multipliers, because they are molecules with
one selected input and one selected output which ”propagate along a FAN-OUT
gate”, while multipliers are multiplied by a FAN-OUT gate. Propagators can
serve in fact as labels, or names, which propagate along a tree of FAN-OUT
gates.
Let’s see, as an application, how the IFTHENELSE construct appear as a
combinator molecule in the chemical concrete machine. In lambda calculus,
there are terms called TRUE, FALSE and IFTHENELSE, which are the of the
booleans true, false and if-then-else. The associated graphs in the chemical
concrete machine are:
Take two molecules A, B, with one exit each. Then IFTHENELSE TRUE A B should
become A. In the chemical concrete machine, with only $\beta^{+}$ enzymes, the
following chain of reactions happens:
Along this chain of reactions, there is no other choice than the ones from the
figure, because at every step there is only one reaction site available to the
enzyme $\beta^{+}$. The result is, unsurprisingly, compatible with the lambda
calculus version, with the exception that A and B are not supposed to be
(graphs corresponding to) lambda terms. They can be anything, for example,
from the family of ”other molecules”.
In lambda calculus IFTHENELSE FALSE A B should become (by reductions) B. In
the chemical concrete machine the following chain of reactions happens:
With a little bit of imagination, if we look closer to what TRUE, FALSE and
IFTHENELSE are doing, we see that it is possible to adapt the IFTHENELSE to a
molecule which releases, under the detection of one molecule (like TRUE), the
”medicine” A, and under the detection of another molecule (like FALSE) the
”medicine” B.
## References
* [1] N. Arkani-Hamed, J. L. Bourjaily, F. Cachazo, A. B. Goncharov, A. Postnikov, J. Trnka, Scattering Amplitudes and the Positive Grassmannian, arXiv:1212.5605
* [2] G. Berry, G. Boudol, The Chemical Abstract Machine, Theoretical Computer Science, 96(1):217–248, 1992
* [3] , M. Buliga, Graphic lambda calculus, to appear in Complex Systems, arXiv:1305.5786
* [4] M. Buliga, Dilatation structures I. Fundamentals, J. Gen. Lie Theory Appl., 1 (2007), 2, 65-95. arXiv:math/0608536
* [5] M. Buliga, Emergent algebras, arXiv:0907.1520
* [6] W.Fontana and L.W.Buss, The Barrier of Objects: From Dynamical Systems to Bounded Organizations, in: Boundaries and Barriers, J.Casti and A.Karlqvist (eds.), pp.56–116, Addison-Wesley, 1996 (pdf)
* [7] W.Fontana and L.W.Buss, ’The Arrival of the Fittest’: Toward a Theory of Biological Organization, Bull. Math. Biol., 56, 1-64 (1994) (pdf)
* [8] W.Fontana and L.W.Buss , What would be conserved if ‘the tape were played twice’?, Proc. Natl. Acad. Sci. USA, 91, 757–761 (1994) (pdf)
* [9] A. Kissinger, Pictures of Processes: Automated Graph Rewriting for Monoidal Categories and Applications to Quantum Computing, arXiv:1203.0202
* [10] G. Moe-Behrens, The biological microprocessor, or how to build a computer with biological parts, Computational and Structural Biotechnology Journal, 7, 8 (Apr 2013) e201304003
* [11] D.P. Thurston, The algebra of knotted trivalent graphs and Turaev’s shadow world, Geom. Topol. Monogr. 4 (2002) 337-362, arXiv:math/0311458
* [12] Varela, F., H. R. Maturana and R. Uribe. 1974. Autopoiesis: the organization of living systems, its characterization and a model, BioSystems 5, 187-196.
|
arxiv-papers
| 2013-09-26T14:24:23 |
2024-09-04T02:49:51.572031
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Marius Buliga",
"submitter": "Marius Buliga",
"url": "https://arxiv.org/abs/1309.6914"
}
|
1309.6923
|
CERN–2013–00512 August 2013Fifty years of the CERN Proton SynchrotronVolume
IIEditors:S. GilardoniD. ManglunkiGENEVA2013
ISBN | 978–92–9083–391–8
---|---
ISSN | 0007–8328
DOI | 10.5170/CERN–2013–005
Copyright © CERN, 2013
Creative Commons Attribution 3.0
Knowledge transfer is an integral part of CERN’s mission.
CERN publishes this report Open Access under the Creative Commons Attribution
3.0 license (http://creativecommons.org/licenses/by/3.0/) in order to permit
its wide dissemination and use.
This monograph should be cited as:
Fifty years of the CERN Proton Synchrotron, volume II
edited by S. Gilardoni and D. Manglunki, CERN-2013-005 (CERN, Geneva, 2013),
DOI: 10.5170/CERN–2013–005
Dedication
The editors would like to express their gratitude to Dieter Möhl, who passed
away during the preparatory phase of this volume. This report is dedicated to
him and to all the colleagues who, like him, contributed in the past with
their cleverness, ingenuity, dedication and passion to the design and
development of the CERN accelerators.
Abstract
This report sums up in two volumes the first 50 years of operation of the CERN
Proton Synchrotron. After an introduction on the genesis of the machine, and a
description of its magnet and powering systems, the first volume focuses on
some of the many innovations in accelerator physics and instrumentation that
it has pioneered, such as transition crossing, RF gymnastics, extractions,
phase space tomography, or transverse emittance measurement by wire scanners.
The second volume describes the other machines in the PS complex: the proton
linear accelerators, the PS Booster, the LEP pre-injector, the heavy-ion linac
and accumulator, and the antiproton rings.
Preface
It was on 24 November 1959 that the proton beam in the CERN Proton Synchrotron
was accelerated to a kinetic energy of 24 GeV. Thus the first strong-focusing
proton synchrotron ever built has been faithfully serving the international
physics community for 50 years. It has been the subject of a virtually
continuous upgrade boosting its intensity per pulse from $10^{10}$ protons by
more than three orders of magnitude to $3\times 10^{13}$ protons. Various
injectors have been added and it has been modified such that, in addition to
protons, light and heavy ions, positrons and electrons, as well as antiprotons
could be accelerated or even decelerated often within the same supercycle.
This would not have been possible had the initial design not been solid and
sound allowing for maintainability, flexibility, and versatility and whose
intrinsic potential was brought to fruition by the efforts and the ingenuity
of generations of accelerator physicists, engineers, operators, and
technicians. This report has been written to mark the fiftieth anniversary of
the first operation of this unique accelerator. Volume I outlines the euphoric
spirit in the European physics community in which such a bold design could be
suggested, and gives an overview of the evolution of this unique accelerator
described in a wealth of publications. This volume provides also a description
in more depth of the outstanding achievements and highlights in its
development. Volume II provides an overview of the injectors of the PS and of
the accelerator system used for antiproton accumulation and storage, which has
been closely associated with the PS.
Contributors
Jean-Paul Burnet, Christian Carli, Michel Chanel, Roland Garoby, Simone
Gilardoni, Massimo Giovannozzi, Steven Hancock, Helmut Haseroth, Kurt Hübner,
Detlef Küchler, Julian Lewis, Alessandra Lombardi, Django Manglunki, Michel
Martini, Stephan Maury, Elias Métral, Dieter Möhl, Günther Plass, Louis
Rinolfi, Richard Scrivens, Rende Steerenberg, Charles Steinbach, Maurizio
Vretenar, Thomas Zickler
Contents
Preface id1
List of contributors id1
Linac 1 LABEL:S01-Linac1
The Proton Synchrotron Booster (PSB) LABEL:S02-PSB
Linac 2 LABEL:S03-Linac2
The Antiproton Accumulator, Collector and Decelerator Rings LABEL:S05-AA-AC-AD
The Low-Energy Antiproton and Ion Rings LEAR and LEIR LABEL:S06-LEAR-LEIR
The LEP Pre-injector (LPI) LABEL:S04-LPI
Linac 3 LABEL:S07-Linac3
See pages 1 of 01-Linac1.pdf See pages 2- of 01-Linac1.pdf See pages 1 of
02-PSB.pdf See pages 2- of 02-PSB.pdf See pages 1 of 03-Linac2.pdf See pages
2- of 03-Linac2.pdf See pages 1 of 05-AA-AC-AD.pdf See pages 2- of 05-AA-AC-
AD.pdf See pages 1 of 06-LEAR-LEIR.pdf See pages 2- of 06-LEAR-LEIR.pdf See
pages 1 of 04-LPI.pdf See pages 2- of 04-LPI.pdf See pages 1 of 07-Linac3.pdf
See pages 2- of 07-Linac3.pdf
|
arxiv-papers
| 2013-09-26T14:44:55 |
2024-09-04T02:49:51.579618
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Simone Gilardoni (ed.) (CERN), Django Manglunki (ed.) (CERN),\n Jean-Paul Burnet, Christian Carli, Michel Chanel, Roland Garoby, Massimo\n Giovannozzi, Steven Hancock, Helmut Haseroth, Kurt H\\\"ubner, Detlef\n K\\\"uchler, Julian Lewis, Alessandra Lombardi, Michel Martini, Stephan Maury,\n Elias M\\'etral, Dieter M\\\"ohl, G\\\"unther Plass, Louis Rinolfi, Richard\n Scrivens, Rende Steerenberg, Charles Steinbach, Maurizio Vretenar, Thomas\n Zickler",
"submitter": "Scientific Information Service CERN",
"url": "https://arxiv.org/abs/1309.6923"
}
|
1309.7083
|
# Characterization of EASIROC as Front-End for the readout of the SiPM at the
focal plane of the Cherenkov telescope ASTRI
D. Impiombato [email protected] S. Giarrusso
[email protected] T. Mineo [email protected] M. Belluso
S. Billotta G. Bonanno O. Catalano A. Grillo G. La Rosa D. Marano G.
Sottile INAF, Istituto di Astrofisica Spaziale e Fisica cosmica di Palermo,
via U. La Malfa 153, I-90146 Palermo, Italy INAF, Osservatorio Astrofisico di
Catania, via S. Sofia 78, I-95123 Catania, Italy
###### Abstract
The Extended Analogue Silicon Photo-multiplier Integrated Read Out Chip,
EASIROC, is a chip proposed as front-end of the camera at the focal plane of
the imaging Cherenkov ASTRI SST-2M telescope prototype. This paper presents
the results of the measurements performed to characterize EASIROC in order to
evaluate its compliance with the ASTRI SST-2M focal plane requirements. In
particular, we investigated the trigger time walk and the jitter effects as a
function of the pulse amplitude. The EASIROC output signal is found to vary
linearly as a function of the input pulse amplitude with very low level of
electronic noise and cross-talk ($<$1%). Our results show that it is suitable
as front-end chip for the camera prototype, although, specific modifications
are necessary to adopt the device in the final version of the telescope.
###### keywords:
Front-end, ASIC for SiPM, Front-end Electronics for Detector, Readout Analogue
Electronics Circuits, Electronic Detector Readout Concepts, Trigger Concepts
and Systems, Cherenkov Telescope
††journal: Nuclear Physics A
## 1 Introduction
The use of Cherenkov telescopes in gamma-ray astronomy opened a new frontier
in the study of the emission from very energetic sources such as supernovae,
neutron stars and supermassive black-holes. Since the discovery of TeV
emission from the Crab nebula with Whipple [1] in 1989, the number of known
sources emitting in this range of energy has rapidly increased to hundreds in
a relatively small number of years. The Cherenkov Telescope Array (CTA) is a
new generation of ground based gamma-ray telescopes aiming at observing the
sky from a few tens of GeV to beyond 100 TeV. It is an array comprising
telescopes with three different sizes each designed to cover a specific range
of energy [2]. The low energy band (from a few tens of GeV up to a hundred
GeV) will be covered by 24 m mirror diameter telescopes that observe a Field
of View (FoV) of 4∘-5∘. The range 100 GeV-1 TeV will be observed by 10-12 m
mirror diameter telescopes having a FoV of 6∘-8∘. The highest energy region
will be observed by small size telescopes (4-6 m mirror diameter) having a FoV
of about 10 degrees.
ASTRI (Astrofisica con Specchi a Tecnologia Replicante Italiana; [3]), is a
’Flagship Project’ of the Italian Ministry of Education, University and
Research led by the Italian Institute of Astrophysics (INAF), whose main
object is to build a prototype for the small size CTA telescopes. The
prototype, hereafter named ASTRI SST-2M telescope, will adopt innovative
solutions to explore the 1-100 TeV range of the electromagnetic spectrum: a
wide field dual mirror Schwarzschild-Couder optical system and Silicon Photo-
Multipliers (SiPM) as sensors of the camera at the focal plane.
The telescope design foresees a focal ratio F# = 0.5, an equivalent focal
length of 2150mm, an average effective area of about 6.5m2, and, considering
the size of the focal plane surface, it covers a Field of View (FoV) of
$\sim$9.6∘ in diameter.
The detection of Very High Energy (VHE) gamma photons is possible by observing
the Cherenkov light produced by the relativistic secondary particles in the
air shower generated by the interaction with the Earth atmosphere. This light,
emitted in the ultraviolet band within a cone of about 1.3∘, reaches the
ground in a ” pool ” having a radius of $\sim$120 m. It is very faint and
lasts only a few ns with a duration that depends mainly on the distance
between the shower core and the telescope axis [4].
SiPMs are suitable for the detection of the Cherenkov flashes, as already
demonstrated by the FACT project [5], because they are very fast and sensitive
to the light in the range 300-900 nm. Their drawbacks with respect to the
traditional PMTs (PhotoMultiplier Tubes) are: a very high dark count, after
pulses, optical cross-talk between elementary diodes inside a pixel and a gain
strongly dependent on chip temperature. Nevertheless, for a single pixel the
rate of the dark counts is lower than the rate of the Night Sky Background
(NSB), so that the high instrumental background does not degrade the telescope
sensitivity. Moreover, the effects of optical cross-talk and after pulses,
that may affect the sensitivity, are typically lower than 20%[6] and the gain
can be kept stable with an adequate strategy of temperature control.
In this paper, we present a set of measurements aimed at characterizing the
performance of EASIROC, a device devoted to read the SiPM output, and compare
the results with ASTRI SST-2M focal plane requirements.
## 2 The ASTRI Camera
The camera at the focal plane of the ASTRI SST-2M is based on monolithic SiPMs
Hamamatsu
S11828-3344m111http://www.hamamatsu.com/sp/hpe/HamamatsuNews/HEN111.pdf with
4x4 squared pixels, 3$\times$3mm large, made up of 3600 elementary diodes of
50$\mu$m pitch giving a filling factor of 62%. The physical pixels are
2$\times$2 grouped in logical pixels (6.2$\times$6.2mm) shown with a yellow
square in Fig. 1. Therefore, considering that, at the focal plane, 1∘
corresponds to 37.5mm, the logical pixel has a size of 0.17∘ matching the
optics angular resolution. The detector units are organized in an array of 37
Photon Detection Modules (PDM) with 8$\times$8 logical pixels each (see Fig.
1), for modularity and a fast read out of the focal plane after any trigger.
The trigger threshold is chosen to have a maximum rate of $\sim$300 Hz from
the whole focal plane that ensures a dead time $<$3%. The energy working range
of the the ASTRI SST-2M telescope prototype is 1-100 TeV and, considering the
optics area, the requirements for the maximum number of photoelectrons
detected in one pixel is 1000 with a goal of 2000.
The very short duration of gamma events requires a dedicated FEE (Front-End
Electronics) capable, not only to catch the very fast pulses of Cherenkov
light, but also to provide auto-trigger capability. The device proposed for
the ASTRI SST-2M telescope prototype is the Extended Analogue Silicon Photo-
Multiplier Integrated Read Out Chip (EASIROC; [7]) equipped with 32-channels
each with the capability of measuring charge from 1 to 2000 photoelectrons
assuming a SiPM gain of 106. Two EASIROC devices are then devoted to read a
single PDM.
Figure 1: A PDM of ASTRI SST-2M focal plane composed of 4$\times$4 SiPMs. The
yellow square shows the logical pixels obtained grouping 2$\times$2 physical
pixels
## 3 The EASIROC chip
EASIROC, produced by Omega Micro222http://omega.in2p3.fr, is a 32 channel
fully analogue front-end Application-Specific Integrated Circuit (ASIC)
dedicated to the readout of SiPM detectors. The architecture of its analog
core and the characteristics are shown in Fig. 2 and Table 1, respectively.
In each channel, two separate electronics chains allow High- and Low-gain
(HG/LG) of the signal in order to measure charge from 160 fC up to 320 pC
covering a range up to 2000 photoelectrons for the adopted $10^{6}$ SiPM gain.
The high voltage of each logical pixel can be individually tuned by varying
the bias voltage in a 8-bit DAC from 0 to 4.5V in order to equalize the gains.
Each of the two chains is composed of an adjustable preamplifier followed by a
tunable shaper and a track and hold circuit. A third chain is implemented to
generate a trigger using a fast shaper followed by a discriminator, whose
threshold is set by a 10-bit DAC common to all the 32 channels (see Fig. 2).
The signal shaping time is set to 50 ns according to our needs, while the
trigger shaping time is fixed to 15 ns.
The power consumption is lower than 5 mW/channel and unused features can be
disabled to further save power. All EASIROC main parameters can be programmed
by downloading a configuration table through a slow control serial line. The
processing of the analog signal takes place in the front-end channels of the
device, while the read-out is handled at the internal back-end of the ASIC.
The outputs of all the channels can be readout from multiplexers that, running
in parallel, can sequentially switch the sampled signal of the two chains (LG
and HG).
An evaluation board has been designed and realized by Omega Micro to test the
functional characteristics and performance of this ASIC. It allows an easy
access to the EASIROC output and provides many test points to the Field
Programmable Gate Array (Altera Corporation - Cyclone FPGA Model EP1C6Q240C6N)
dedicated lines (see Fig. 3). It is equipped with two external Analog-to-
Digital Converters (ADC) to allow ASIC digital data acquisition. The trigger
of each channel and the OR output of the 32 trigger channels ($OR32$) are
available in parallel. The 32 triggers are sent to the FPGA that, according to
the set threshold, generates a common hold signal ($HOLD-B$) for the output
read-out. A LabVIEW software of the National Instruments, developed by the LAL
(Laboratoire de l’Acc$\acute{e}$l$\acute{e}$rateur Lin$\acute{e}$aire) Tests
group333 http://www.lal.in2p3.fr/, has been provided together with the
evaluation board to command the EASIROC chip and to receive the output via
Universal Serial Bus (USB) connection. It allows one to send the ASIC
configuration and to receive the output bits via a USB cable connected to the
test board.
In our measurements, we used a new version of the chip, EASIROC-A, kindly
provided by Omega Micro.
Table 1: Main characteristics of the EASIROC chip. Technology: | Austria-Micro-Systems (AMS) SiGe 0.35 $\mu$m
---|---
Dimensions : | 16.6 mm2 (4.15$\times$4.01mm)
Power Supply : | 4.5V/0V
Consumption: | 4.84mW per channel
Package : | Naked (PEBS) TQFP160
Figure 2: Architecture of the front-end EASIROC (Omega Micro courtesy) Figure
3: Architecture of the PCB (Printed Circuit Board) of the EASIROC (Omega Micro
courtesy)
## 4 Laboratory set-up
A pulse function arbitrary generator is used to create the input charge
injected in EASIROC. The signal is similar, as much as possible, to the SiPM
one characterized by very fast rise time (a few hundreds of ps) followed by an
exponential decay a few tens of ns). The amplitude of this signal for an input
charge of 0.16 pC, equivalent to 1 pe for a SiPM gain of 106, is
300$\pm$4$\mu$V.
According to the goal on the ASTRI SST-2M maximum rate per pixel, the gain of
the LG chain preamplifier is set to 3 that produces a monotonic response up to
2000 pe. Moreover, we decided to fix the gain of the HG chain preamplifier to
150 allowing an almost linear working range up to $\sim$50 pe (see Sect. 6.3).
The shaping time of the signal was fixed to 50 ns, considering that the higher
energy events, mainly coming from large shower core distance, are expected to
last up to $\sim$30 ns.
In our sets of measurements, we first characterised the trigger behaviour and
the delays introduced by the jitter and by the time walk. We then sampled the
shape of the output signal for both HG and LG chains and investigated its
linearity with respect to the injected charge. Finally, we evaluated the level
of the cross-talk. All the measurements were performed using only one of the
32 channels (Ch 31). However, to check the validity of the results, some
measurements were repeated for other channels.
## 5 Trigger
The trigger chain, derived from the HG preamplifier, is composed of a
dedicated 15 ns fast shaper followed by a discriminator which provides the
trigger signal. The threshold of the discriminator is common to all the 32
channels and it is set by a 10-bit DAC whose output range goes from 1.06V up
to 2.38V, in steps of 1.28mV.
To characterise the DAC-controlled discriminator, we tested the threshold
voltage at the 1024 DAC values using a Keithley 2000 Digital Multimeter and we
found that Voltages follow a straight line with an average discrepancy of
$\sim$0.5$\%$.
The trigger output was evaluated varying the discriminator threshold for fixed
injected charge in the range 0-8 pC (0-50 pe) using a reference clock (10kHz)
synchronous to the signal. As an example, we show in Fig. 4 the efficiency
curves for 0.33 pe, 1 pe, 2 pe and 3 pe. The last dotted line at $\sim$930 DAC
corresponds to the trigger efficiency with no injection charge, obtained
keeping the same reference clock. In this case, triggers are due to the
electronics noise.
From simulations, we evaluated that the requirement of 300Hz on the total
trigger rate at the focal plane is ensured if at least four contiguous pixels
are triggered, with the threshold set at 3 pe. The trigger efficiency as
function of the injected charge is shown in Fig. 5 for a threshold set to 3
pe.
Figure 4: Trigger efficiency for fixed injected charges as function of the
discriminator threshold. The dotted line at 930 DAC counts corresponds to the
trigger efficiency with no injection charge Figure 5: Trigger efficiency vs
injected charge setting the discriminator threshold to 3 pe.
### 5.1 Time Walk and Jitter
The time walk and jitter of the trigger were investigated by measuring, at the
oscilloscope, the delay between the trigger at the chip and the strobe of the
generator. The range of injected charge is 0.16–3.2 pC (1-20 pe) and the
discriminator threshold level is sequentially set to 0 pe (the pedestal), 1
pe, 2 pe and 3 pe.
Since these measurements could be influenced by the jitter in the signal, we
first investigated its level and found that it is negligible being of the
order of a few tens of ps.
The time jitter, measured as the RMS of the delay distribution, are presented
in Fig. 6 as a function of injected charge. As expected, it is higher at
pedestal but, in the sampled charge range, its value is always lower than 1.23
ns and decreases with increasing threshold. For a discriminator threshold of 3
pe and injected charges higher than 0.64 pC (4 pe) it is lower than 0.3 ns.
We show in Fig. 7 the values of the time walk as function of injected charges.
We measured them as the mean of the delay distribution and shifted them so
that the time walk corresponding to the largest injected charge is zero in
each case. The time walk is always below 5.5 ns and it is reduced to 1 ns for
signals higher than 1.92 pC (12 pe) for a discriminator threshold of 3 pe.
Figure 6: Jitter of the trigger vs injected charge at several discriminator
thresholds. Figure 7: Time walk of the trigger vs injected charge at several
discriminator thresholds. The zero level was set to the time walk
corresponding to the largest injected charge.
## 6 Output signal
The expected range of the Cherenkov signal (1-2000 pe) requires the
characterization of both HG and LG chain. Our measurements investigated the
level of the electronic noise for different working conditions of the chip,
the shape of the output signal and the linearity as function of the injected
charge.
### 6.1 Electronic Noise
The level of the noise affecting the signal was evaluated as function of the
preamplifier gain. We defined the noise as the $RMS$ of the analogue output
signal with no charge injected. To convert the voltages into Equivalent Noise
Charge (ENC), we used the following formula:
$ENC\,=\,K\times V_{RMS}\hskip 14.22636pt{\rm with}\hskip
14.22636ptK=\frac{Q_{Npe}}{V_{Npe}}$ (1)
where $V_{RMS}$ is the rms spread in output voltage due to the noise,
${Q_{Npe}}$ and ${V_{Npe}}$ are the charge and the output voltage
corresponding to a signal $Npe$, respectively. We found that $K=0.045$ [pe
mV-1] for gain = 150 in the HG and $K=1.15$ [pe mV-1] for gain = 3 in the LG
chain.
The feedback capacitance (see Fig. 2) was varied in its working range 0.1-1.5
pF corresponding to the gain range 150-10 in the HG chain and to 15–1 in the
LG one. The average noise together with its variability range are shown in
Fig. 8, for both gains. In the HG chain, the noise is 2.1 mV at gain =150,
which converts to 0.1 pe using the Equation 1.
The noise in the LG chain is constant for all the investigated gains and its
average value is 0.91$\pm$0.05 mV, equivalent to $\sim$1 pe.
Figure 8: Average electronic noise in mV measured at the oscilloscope vs the
preamplifier gain. The vertical dashed lines indicate the gains at which the
EASIROC will be operated. The shaded areas represent the noise variability
ranges. The left panel relates to the HG and the right panel to the LG.
### 6.2 Slow-shaper sampling
The chip saves the amplitude of the preamplified and shaped signal at a given
time, the hold time, directly set by the user using a track and hold cell (see
Fig. 2). Two external ADCs, one for the HG and one for the LG, digitize the
signal level measured at the hold time (see Fig. 2).
We sampled the shaped signal for the two electronics chains at several hold
times with the main task of evaluating the peaking time. Each point was
derived averaging over 1000 samplings and several different injected charges
were investigated fixing the shaping time to 50 ns.
The pedestal-subtracted signal was normalized to give a peak amplitude of one,
and the time axis was adjusted so that this peak was centred at 0. The curves
as a function of time in ns are shown in 9 and 10 for the HG and LG chains,
respectively. We note that, 10 ns after the peaking time, the amplitude of the
signal produced by an injected charge of 5.6 pC is reduced by $\sim$25 ADC
counts with respect to the peak (corresponding to $\sim$1.4% in the normalized
curve of Fig. 9) for the HG chain, and by $\sim$1 ADC counts (corresponding to
$\sim$1% in the curve of Fig. 10) for the LG chain.
In order to compute the peaking time, the peak of the signal was fitted with a
Gaussian whose best fit values are presented in Fig. 11 for the two gain
chains. We note that, the peaking time of the LG is quite constant: it varies
by 0.6 ns with respect to its average (47.1 ns) over the entire investigated
range (56-338 pe). The case is different in the 6-50 pe range of the HG chain
where the peaking time varies by about 6% with respect to the average of 74.0
ns that is moreover 24 ns higher than the designed shaping time (50 ns).
Figure 9: Amplitude normalized to the ADC value at the peaking time versus
the time in ns shifted at the peaking time for the HG chain. Figure 10:
Amplitude normalized to the ADC value at the peaking time versus the time in
ns shifted at the peaking time for the LG chain.
Figure 11: Peaking time vs the injected charge. Red points are relate to the
LG chain, blue points to the HG one.
### 6.3 Charge linearity
The linearity of the system was checked by analyzing the slow shaper amplitude
at the time 50 ns for the LG chain and 80 ns for the HG one as function of the
injected charge. We assumed as error of each point the sigma of the
distributions obtained performing 1000 tests. These uncertainties are always
lower than 1%. Furthermore, the Integral and Differential Nonlinearity
specified for the device AD9220 is much smaller than this
uncertainty444http://www.alldatasheet.com/datasheet-
pdf/pdf/237229/AD/AD9220.html. The ADC values as function of injected charges
are shown in Fig. 12 and Fig. 13 for the LG and the HG.
We note that the ADC values curve has a monotonic behaviour up to 8 pC (50 pe)
for the HG and up to 320 pC (2000 pe) for the LG. These two working ranges are
then suitable for the ASTRI SST-2M requirements (see Sect.4).
Moreover, we established the range where the correlation between the ADC
values and the charge can be modelled with a straight line. Starting from a
few points we increased the fitting range up to the maximum charge that still
gives an acceptable $\chi^{2}$. We find that the HG chain is linear in the
range 0.16-6.4 pC and the LG chain in the range 0.16-64 pC, where $\chi^{2}$
=6.5 for 6 degrees of freedom(dof) and $\chi^{2}$ =13 (7 dof) were obtained
for HG and LG, respectively. The measured points and the fitting lines in the
validity ranges are shown in Fig. 12 and Fig. 13.
Figure 12: ADC values vs injected charges for the HG chain. The continuous
line is the linear best fit up to the injected charge of 6.4 pC (see text).
Figure 13: ADC values vs injected charges for the LG chain. The continuous
line is the linear best fit up to the injected charge of 64 pC (see text).
## 7 Cross-Talk
Our last step of measurements were devoted to investigate the cross-talk as
function of the injected charge. This was performed measuring the output
signal in the two $Ch\leavevmode\nobreak\ 31$ neighboring channels
($Ch\leavevmode\nobreak\ 0$ and $Ch\leavevmode\nobreak\ 30$) for input charges
ranging from 1 to 320 pC (6-2000 pe). We assumed that a cross-talk signal is
detected if its intensity is 2.5 $\sigma$ above the noise. We find that no
significant cross-talk was observed for the LG in the whole investigated range
and for the HG below 5 pC (31 pe).
The ratios, in percentage, between the outputs in $Ch\leavevmode\nobreak\ 0$
and $Ch\leavevmode\nobreak\ 30$ and the output in $Ch\leavevmode\nobreak\ 31$
for the HG chain for input charges above 5 pC are presented in Fig. 14. We
note that the effect of the cross-talk is always lower than 1% in the ASTRI
SST-2M HG working range ($<$50 pe) and its level is $<$3% in the range where
HG saturates (bottom panel of Fig. 14).
These values must be compared with the cross-talk between logical pixels,
which is not known at the moment. Dedicated measurements are foreseen in the
SiPM characterization.
Figure 14: The ratios, in percentage, between the outputs in
$Ch\leavevmode\nobreak\ 0$ and $Ch\leavevmode\nobreak\ 30$ and the output in
$Ch\leavevmode\nobreak\ 31$ vs injected charge for the HG chain. The top panel
corresponds to the ASTRI working range ($<$50 pe); the bottom panel shows the
cross talk in the range where the HG chain saturates.
## 8 Summary and Conclusion
We performed a set of measurements to investigate the performance of EASIROC
as the ASIC front-end for the camera at the focal plane of the ASTRI SST-2M
telescope prototype. The measurements were mainly aimed at understanding the
triggering capability and the shape of the output signals.
Table 2: Experimental results of EASIROC characterization. Jitter(discriminator thresholds 3 pe) | 0.3 ns
---|---
Time walk(discriminator thresholds 3 pe) | $<$5.5ns
Noise HG | $\sim$0.1 pe
Noise LG | $\sim$1 pe
Cross Talk HG | $<$1$\%$
Cross Talk LG | No
The main experimental results are presented in the Table 2. They are
summarized in the following points and compared with the ASTRI SST-2M
requirements:
* 1.
Trigger: EASIROC is able to trigger events with a number of counts $>$0.3 pe,
well below the value foreseen for the ASTRI SST-2M prototype. In fact, the
requirement on the maximum trigger rate (300 Hz) can be obtained if at least
four contiguous pixels each have a signal greater than 3 pe.
The trigger time walk and jitter are always below 5.5 ns and 0.3 ns,
respectively, for a trigger threshold set to 3 pe. They introduce an
uncertainty in the charge measurements of a few percent in both chains, with
the shaping function of EASIROC configured to give a signal with a broad peak.
* 2.
Signal: The output signal has a stable peaking time in the LG chain and its
value corresponds to the set shaping time. Some problems are detected in the
HG chain. The peaking time is on average shifted by $\sim$20 ns with respect
to the shaping time with a significant drift as function of the injected
charge. Since the hold signal $HOLD-B$ is common to both the chains in the
present version of the chip, a reduction in the amplitude of the HG signal is
measured (see Fig.9). This effect can be reduced, without degrading the
performance of the LG chain, by choosing a suitable sampling time. However, a
new version of the chip should correct this shortcoming either by anticipating
the HG peaking time or by providing two different hold signals.
We detected a monotonic behaviour up to 8 pC for HG chain and up to 320 pC for
the LG one. The levels of the electronic noise and the cross-talk between
channels introduced by EASIROC are negligible.
In conclusion, our measurements showed that EASIROC is suitable for the front-
end of the camera at the focal plane of the ASTRI SST-2M telescope prototype.
Minor ad-hoc modifications in its re-design, as for example a different $HOLD$
signal for each electronics chain, are however necessary to tailor this device
to be fully compliant with the telescope requirements.
## Acknowledgements
The work presented in this paper was partially supported by the ASTRI,
”Flagship Project” financed by the Italian Ministry of Education, University,
and Research (MIUR) and lead by the Italian National Institute of Astrophysics
(INAF). We are deeply grateful to S. Callier, C. De La Taille, and L. Raux of
the Omega Micro at Orsay and to M.C. Maccarone, A. Rubini and G. Tosti of the
ASTRI collaboration for useful discussions and suggestions.
## References
* [1] T.C. Weekes, et al., ”Observation of TeV gamma rays from the Crab nebula using the atmospheric Cerenkov imaging technique”, ApJ, 342, p.379-395, 1989.
* [2] M. Actis et al., The CTA Consortium, ”Design concepts for the Cherenkov Telescope Array CTA: an advanced facility for ground-based high-energy gamma-ray astronomy”, Experimental Astronomy, Vol.32, Issue 3, p. 193-316, 2011.
* [3] R. Canestrari, O. Catalano, P. Conconi, G. Pareschi, B. Sacco, S. Vercellone, on behalf of the ASTRI collaboration, ”The Italian ASTRI program: an end-to-end dual-mirror telescope prototype for Cherenkov light imaging above few TeV”, 32nd International Cosmic Ray Conference, Beijing, 2011.
* [4] M. He$\ss$, K. Bernlöhr, A. Daum, et al., ”The time structure of Cherenkov images generated by TeV $\gamma$-rays and by cosmic rays”, Astroparticle Physics 11, 363-377, 1999.
* Anderhub et al. [2009] H. Anderhub, M. Backes, A. Biland, et al., ”A novel camera type for very high energy gamma-ray astronomy based on Geiger-mode avalanche photodiodes”, Journal of Instrumentation, 4, 10010, 2009.
* [6] P. Eckert, H.C. Schultz-Coulon, W. Shen, R. Stamen, A. Tadday, ”Characterisation Studies of Silicon Photomultipliers”, Nuclear Instruments and Methods in Physics Research Section A, Vol. 620, Issues 2 3, p. 217-226, 2010.
* [7] S. Callier, C. de La Taille, G. Martin-Chassarda, L. Raux, ”EASIROC, an easy & versatile ReadOut device for SiPM”, TIPP 2011 Technology and Instrumentation in Particle Physics 2011.
* [8] D. Impiombato, S. Giarrusso, et al., ”Characterization of the front-end EASIROC for read-out of SiPM in the ASTRI camera”, Procs. SciNeGHE 2012 Workshop, 20-22 June 2012, Lecce, Italy, Nuclear Physics B: Proceedings Supplements, in press.
|
arxiv-papers
| 2013-09-26T22:34:11 |
2024-09-04T02:49:51.588973
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "D. Impiombato, S. Giarrusso, T. Mineo, M. Belluso, S. Billotta, G.\n Bonanno, O. Catalano, A. Grillo, G. La Rosa, D. Marano, G. Sottile",
"submitter": "Domenico Impiombato",
"url": "https://arxiv.org/abs/1309.7083"
}
|
1309.7148
|
arxiv-papers
| 2013-09-27T08:30:12 |
2024-09-04T02:49:51.597744
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Guo-Long Tan, Li-Min Liu, Weibing Wu",
"submitter": "Guolong Tan",
"url": "https://arxiv.org/abs/1309.7148"
}
|
|
1309.7352
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# Evolution of star clusters in a cosmological tidal field
Steven Rieder1,2, Tomoaki Ishiyama3, Paul Langelaan1, Junichiro Makino4,
Stephen L. W. McMillan5 and Simon Portegies Zwart1
1 Sterrewacht Leiden, Leiden University, P.O. Box 9513, 2300 RA Leiden, The
Netherlands
2 Section System and Network Engineering, University of Amsterdam, Amsterdam,
The Netherlands
3 Center for Computational Sciences, University of Tsukuba, Japan
4 Interactive Research Center of Science, Graduate School of Science and
Engineering Tokyo Institute of Technology,
2-12-1 Ookayama, Meguro, Tokyo 152-8551, Japan
5 Department of Physics, Drexel University, Philadelphia, PA 19104, USA
[email protected]
(Accepted …. Received …; in original form …)
###### Abstract
We present a method to couple $N$-body star cluster simulations to a
cosmological tidal field, using the Astrophysical Multipurpose Software
Environment. We apply this method to star clusters embedded in the CosmoGrid
dark matter-only $\Lambda$CDM simulation. Our star clusters are born at
$z=10$ (corresponding to an age of the Universe of about 500 Myr) by selecting
a dark matter particle and initializing a star cluster with 32,000 stars on
its location. We then follow the dynamical evolution of the star cluster
within the cosmological environment.
We compare the evolution of star clusters in two Milky-Way size haloes with a
different accretion history. The mass loss of the star clusters is continuous
irrespective of the tidal history of the host halo, but major merger events
tend to increase the rate of mass loss. From the selected two dark matter
haloes, the halo that experienced the larger number of mergers tends to drive
a smaller mass loss rate from the embedded star clusters, even though the
final masses of both haloes are similar. We identify two families of star
clusters: native clusters, which become part of the main halo before its final
major merger event, and the immigrant clusters, which are accreted upon or
after this event; native clusters tend to evaporate more quickly than
immigrant clusters. Accounting for the evolution of the dark matter halo
causes immigrant star clusters to retain more mass than when the z=0 tidal
field is taken as a static potential. The reason for this is the weaker tidal
field experienced by immigrant star clusters before merging with the larger
dark matter halo.
###### keywords:
galaxies: star clusters, galaxies: evolution, cosmology: dark matter, methods:
N-body simulations
††pagerange: Evolution of star clusters in a cosmological tidal
field–LABEL:lastpage
## 1 Introduction
Globular clusters are dense self gravitating systems of a few $10^{4}$ to
$\sim 10^{6}$ stars (Brodie & Strader, 2006). With an age of about 12.6 Gyr
(Krauss & Chaboyer, 2003) they are among the oldest objects in the universe
and tend to populate the haloes of galaxies. Their age is consistent with
being born between $z=12$ to 7, which is consistent with the results of
$\Lambda$CDM simulations (Kravtsov & Gnedin, 2005).
From the time the clusters were born on galaxies grow by about a factor 100 in
mass via mergers to their current mass, size and morphology (White & Rees,
1978; Peacock, 1999). The environment in which the globular clusters evolved
since their birth has consequently changed quite dramatically over their
lifetimes. These changes may have a profound effect on the evolution of star
clusters.
Most modern star-cluster simulations take some sort of background potential of
the host galaxy into account. This started already in the early 1990s with
Chernoff & Weinberg (1990), and soon afterwards became a lively industry. Many
simulations have been performed with a fixed tidal limit (Vesperini & Heggie,
1997; Portegies Zwart et al., 1998; Giersz, 2001) whereas other include some
sort of tidal potential with a more fluent description of the tidal field
(Gnedin & Ostriker, 1997; Takahashi & Portegies Zwart, 2000; Lamers, Baumgardt
& Gieles, 2010). In most of these simulations the cluster orbit was circular
and did not change with time. In a few cases the orbit was allowed to be
eccentric, but still did not change with time (Baumgardt & Makino, 2003;
Giersz & Heggie, 2011, 2009). The next refinement was the relaxing of the
orbital parameters, allowing the cluster orbit to change (Tanikawa &
Fukushige, 2005, 2010; Praagman, Hurley & Power, 2010; Renaud, Gieles & Boily,
2011; Berentzen & Athanassoula, 2012; Renaud & Gieles, 2013). The evolution of
star clusters in a live galactic potential combined with parametrized cluster
evolution was studied by Kruijssen et al. (2011), and Matsui et al. (2012)
performed simulations in which they resolve the formation of star clusters in
a single galaxy merger event.
However, the mass evolution of the parent galaxy is generally ignored in all
these simulations, except for the few cases with a galaxy merger (Saitoh et
al., 2009; Renaud, Gieles & Boily, 2011; Kruijssen et al., 2012; Renaud &
Gieles, 2013).
In this paper we study the evolution of star clusters in a cosmological
environment. The background potential against which the star clusters are
evolved, are taken from the CosmoGrid dark-matter only $\Lambda$CDM
simulation (Portegies Zwart et al., 2010). We selected two Milky Way like
haloes in which the star clusters are simulated. The coupling between the
cosmological simulation and the star cluster is realized via the Astrophysical
Multipurpose Software Environment (AMUSE) (Portegies Zwart et al., 2009, 2011,
2013; Pelupessy et al., 2013). In these simulations we initialize a total of
30 star cluster at $z=10$, and evolve them together with the cosmological
simulation up to $z=0$.
## 2 The experimental setup
We simulate star clusters in a $\Lambda$CDM environment. We do this in two
steps, first by calculating the $\Lambda$CDM environment and then using the
tidal field from this environment as an external tidal field for the star
cluster simulations. We investigate the results of two distinct regions in the
$\Lambda$CDM environment.
### 2.1 The CosmoGrid $N$-body simulation
The cosmological simulation employed in this article originates from the
CosmoGrid calculation (Portegies Zwart et al., 2010; Ishiyama et al., 2013),
which is a dark matter-only $\Lambda$CDM simulation of $2048^{3}$ particles in
a $(21~{}h^{-1}{\rm Mpc})^{3}$ co-moving cosmological volume. We performed
these simulations using the GreeM (Ishiyama, Fukushige & Makino, 2009;
Ishiyama, Nitadori & Makino, 2012) and Sushi (Groen et al., 2011) codes. GreeM
is a massively parallel TreePM code based on the implementation of Yoshikawa &
Fukushige (2005). The SUSHI code is an extension of the GreeM code, which can
run on a planet wide grid of supercomputers. Within both codes, the equations
of motion are integrated in co-moving coordinates using the leap-frog scheme
with a shared, adaptive timestep. In this simulation each particle has a mass
of $1.28\times 10^{5}$ ${\rm M}_{\odot}$. In total, we have 556 snapshots,
separated by $\mbox{dt}\simeq 35\mbox{Myr}$ (for $t<7.5\mbox{Gyr}$) and
$\mbox{dt}\simeq 17.5\mbox{Myr}$ (for $t\geq 7.5\mbox{Gyr}$). The CosmoGrid
simulation lasted from $z=65$ to $z=0$. We employed the following cosmological
parameters: $\Omega_{0}=0.3,\lambda_{0}=0.7,h=0.7,\sigma_{8}=0.8,n=1.0$. For
more details on the simulation see Portegies Zwart et al. (2010); Ishiyama et
al. (2013).
### 2.2 Halo catalogue
We use the halo finder Rockstar (Behroozi, Wechsler & Wu, 2013) to identify
haloes in each snapshot. Rockstar is based on adaptive hierarchical refinement
of friends-of-friends groups in six dimensions and allows for the robust
tracking of subhaloes. We use the gravitationally consistent merger tree code
from Behroozi et al. (2013) to construct the merger history for all haloes
identified by Rockstar. We use AMUSE to find the radial density profiles for
our haloes.
From the $z=0$ haloes, we then select two haloes, based on their relative
isolation and a mass comparable to that of the Milky-Way Galaxy. The two
haloes are quite similar in many respects, but have a different merger history
and the number of subhaloes at $z=0$ is different.
In Fig. 1, we present the merger history of the two haloes, both schematically
and visually. Halo A completes a major merger at around $t=6.5$ Gyr, the halo
it merges with can be seen in the third figure from the top. At the end of the
simulation (see Fig. 2 a), it is in the process of merging with another
similar-sized halo. Other than that, there is no significant interaction
between $t=6.5$ Gyr and $t=13.7$ Gyr. By the end of the simulation, halo A is
the largest object within a radius of $6.3~{}h^{-1}{\rm Mpc}$ (see Tab. 1).
Halo B exists in a denser part of the volume, and as a result more structure
is visible in the outskirts of its environment (see Fig. 2b). By $z=0$, it is
the largest halo within only $0.97~{}h^{-1}{\rm Mpc}$. During its history, it
underwent many small merger events, and one long-lasting major merger event
that completed around $t=11$ Gyr.
In the mass evolution of both haloes (Fig. 3), the larger merger events are
clearly visible. Since the virial mass includes the mass from subhaloes, the
mergers are visible here at the start of interaction, rather than at the end
as in Fig. 1.
The halo density profiles at $z=0$ (see Fig. 4) are consistent with the haloes
described in Ishiyama et al. (2013), with concentration parameters $c_{\rm
vmax}$ of $3.53$ and $3.76$ for haloes A and B respectively. The lower
concentration of halo A may be explained by its ongoing major merger event,
which causes the halo to have two cores (see Fig. 2a).
Table 1: Properties of the selected haloes in the final CosmoGrid snapshot.
Halo | $M_{\rm vir}$ | $R_{\rm vir}$ | $V_{\rm max}$ | $c_{\rm vmax}$ | $N_{\rm sub}$ | $D_{\rm n}$ | b/a | c/a
---|---|---|---|---|---|---|---|---
| $10^{11}$ $h^{-1}{\rm M}_{\odot}$ | $h^{-1}{\rm kpc}$ | km/s | | | $h^{-1}{\rm Mpc}$ | |
A | $6.33$ | $173.8$ | 140.2 | 3.53 | 61 | 6.30 | 0.81 | 0.63
B | $4.78$ | $159.4$ | 133.1 | 3.76 | 29 | 0.97 | 0.73 | 0.68
$M_{\rm vir}$ and $R_{\rm vir}$ are the virial mass and radius (Bryan &
Norman, 1998), $V_{\rm max}$ is the maximum of the rotation curve, $c_{\rm
vmax}$ is the concentration parameter, $N_{\rm sub}$ is the number of
subhaloes with a mass larger than $10^{8}$ $h^{-1}{\rm M}_{\odot}$, $D_{\rm
n}$ is the distance of the nearest more massive halo. a, b and c are the
principal axes of the halo.
Figure 1: The merger history of the two selected dark-matter haloes A (left)
and B (right). The size of each circle is proportional to the virial mass of
the halo. Only haloes and subhaloes with a peak mass larger than $5\times
10^{8}$$h^{-1}{\rm M}_{\odot}$that are accreted at $z=0$ are plotted. The
bottom two halo images are identical to those in Fig. 2.
(a)
(b)
Figure 2: Projected image of the two selected haloes A (top panel) and B
(bottom panel), in the final snapshot ($z=0$). The linear dimension of the
image is $400~{}h^{-1}{\rm kpc}$. The intensity represents the column density
(scaled to minimum/maximum values) and the colour is scaled to the velocity
dispersion. Each galaxy halo contains about 5 million dark-matter particles.
Figure 3: Evolution of the virial mass of the two selected haloes A and B.
Halo A experiences major merger events around 5 Gyr and 11 Gyr, while halo B a
long-lasting major merger event from around 8 Gyr on. Figure 4: Radial
density profiles of the two selected haloes A and B and their environment out
to $400~{}h^{-1}{\rm kpc}$ at $z=0$. The clusters are selected at
galactocentric radii of $3$, $6$ and $12~{}h^{-1}{\rm kpc}$.
### 2.3 The clusters
We select a total of 30 dark-matter particles from the $z=0$ snapshot. For
both haloes we select 15 particles at random, equally divided over three bins
at galactocentric radius $3\pm 0.05~{}h^{-1}{\rm kpc}$, $6\pm
0.05~{}h^{-1}{\rm kpc}$ and $12\pm 0.05~{}h^{-1}{\rm kpc}$. These selected
particles are considered the globular clusters for which we will calculate the
evolution. Since we do not apply further restrictions in the selection
criteria, the clusters may (and likely will) have their peri- and apocentres
well outside these bins. In Fig. 5, we show a projection of these particles in
the central region of their host halo at $z=0$.
(a)
(b)
Figure 5: Projected image of the central portion of the two selected haloes A
(top panel) and B (bottom panel) with their ‘globular clusters’ at $z=0$. The
image size is $24~{}h^{-1}{\rm kpc}$. Red squares indicate clusters at
$3~{}h^{-1}{\rm kpc}$ from the halo centre, green circles those at
$6~{}h^{-1}{\rm kpc}$ and blue triangles those at $12~{}h^{-1}{\rm kpc}$.
We locate the selected dark-matter particles at $z=10$, which corresponds to
an age of the universe of $\sim 0.5$ Gyr. In Fig. 6 we present the $z=10$
image of the two selected haloes with their selected dark-matter particles
that will represent globular clusters, and the distributions from which these
particles are drawn. In both haloes, the particles that end up in the more
central parts of the halo are already the most concentrated in density centres
at $z=10$. The particles of halo A are largely concentrated in two regions, in
which the haloes that merge around $t=6.5$ Gyr form. In contrast, the
particles of halo B are more spread out, reflecting the more violent history
of this halo.
(a)
(b)
Figure 6: Projected image of halo A (top) and B (bottom) at $z=9.65$. The
image size is $2~{}h^{-1}{\rm Mpc}$. Red squares indicate clusters that at
$z=0$ were selected at $3~{}h^{-1}{\rm kpc}$ from the halo centre, green
circles those at $6~{}h^{-1}{\rm kpc}$ and blue triangles those at
$12~{}h^{-1}{\rm kpc}$. The panels on the right display (from top to bottom)
the distributions of particles that end up in spherical shells at $3$, $6$ and
$12~{}h^{-1}{\rm kpc}$ from the halo centres at $z=0$.
At the $z=10$ snapshot we initialize a ‘globular cluster’. Our clusters have
32 000 stars distributed in a Plummer (1911) sphere with a virial radius of 3
pc. The clusters are assumed to be born in virial equilibrium. All stars in
the clusters have the same mass and we did not include stellar evolution.
The star clusters are simulated using the AMUSE framework (Portegies Zwart et
al., 2009, 2011, 2013; Pelupessy et al., 2013). Our simulation code solves for
the equations of motion using Bonsai (Bédorf, Gaburov & Portegies Zwart, 2012)
and ph4 (McMillan et al., 2012). Bonsai is a Barnes & Hut (1986) tree code
that runs on GPUs. It supports both shared timesteps and block timesteps, the
latter allowing individual blocks to have different timesteps for increased
accuracy in dense regions without slowing the simulation down too much. For
this code, we adopted an opening angle (which controls the accuracy, smaller
angles being more accurate) of 0.6, we set the smallest timestep to be used to
1/65536 $N$-body time units (Heggie & Mathieu, 1986) and we used a softening
length of 0.00125 $N$-body length units (0.00375 parsec). Ph4 is a direct
$N$-body integrator with block time steps and GPU acceleration. In order to
directly compare the results to those obtained with Bonsai, we apply the same
softening length as before to the runs with ph4.
In order to validate the use of the tree code, we compare a cluster simulated
with Bonsai to a reference simulation using ph4. In Fig. 7, we present the
mass and the Lagrangian radii of this simulated cluster as a function of time,
for both Bonsai and ph4.
The difference in mass evolution between ph4 and Bonsai remains quite small
until about 5 Gyr. After this moment, both clusters go into core collapse and
the two codes start to deviate more. Until about 8.5 Gyr, the ph4 cluster
displays much higher mass loss than the Bonsai cluster as it expands following
core collapse. After 8.5 Gyr, both codes again show similar behaviour.
The Lagrangian radii of the clusters are nearly equal until core collapse
occurs at about 5 Gyr. After this, the core collapse is initially deeper in
ph4, while after 8.5 Gyr Bonsai reaches the same depth.
From these results, we infer that our Bonsai simulations are not as well
suited for determining the internal structure and evolution of the star
clusters as ph4 would be, and underestimates mass loss due to core collapse.
However, the effect of the tidal field on the mass loss rate is similar in ph4
and Bonsai. Since we investigate only the mass evolution of the clusters due
to the tidal field in this article, we conclude that Bonsai gives an adequate
indication of the effect of tidal fields on the cluster mass loss, and that it
can be used to study the survivability of star clusters.
Figure 7: Evolution of the bound mass (top panel) and the 90%, 50% and 1%
Lagrangian radii (top to bottom, bottom panel) of star cluster A12-2. The
cluster was simulated with Bonsai (red, solid curves) and ph4 (green, dashed
curve).
### 2.4 The tidal field
In each snapshot of the CosmoGrid simulation we calculate the tidal tensor at
the location of the selected dark-matter particle, which represents a star
cluster. The contribution to the tidal tensor $\mathbf{T}_{\rm t}$ from a
particle with relative position $\mathbf{r}$ is given by the second derivative
of the gravitational potential $\phi$:
$T_{t}^{ij}(\mathbf{r^{\prime}})=-\frac{\partial^{2}\phi}{\partial{r^{\prime
i}\partial r^{\prime j}}}$ (1)
where $\mathbf{r}^{\prime}=\mathbf{r}+\epsilon$. For the CosmoGrid simulation,
the value for the softening length $\epsilon$ was 175 parsec.
The strength of the tidal field scales as
$\frac{\partial^{2}\phi}{\partial^{2}\mathbf{r}}\sim\frac{1}{r^{3}}.$ Any
particle at a distance of about $\epsilon$ will have as much effect on the
tidal tensor as the largest halo in the CosmoGrid simulation (containing $\sim
3\times 10^{8}$ particles) would have at a distance of $\sim 117$ kpc. We
therefore include the contribution from all particles within a radius of 125
kpc from our clusters to determine $\mathbf{T}_{\rm t}$.
The strength of the tidal field is calculated from the eigenvalues
$\lambda_{i}$ and eigenvectors $\mathbf{\nu}_{i}$ of this tensor
$\mathbf{T}_{\rm t}$. The eigenvalues give the magnitude of tidal field,
whereas the eigenvectors give the direction along which the system is
stretched. This method for calculating the tidal tensor is similar to the one
employed in Renaud, Gieles & Boily (2011).
Since we calculate the tidal tensor $\mathbf{T}_{\rm t}$ from snapshots of the
CosmoGrid simulation and the number of snapshots is limited, we do not have a
continuous tidal field. In order to prevent large, sudden changes, we linearly
interpolate the tidal tensor between snapshots to create a continuous tidal
tensor. In Sec. 2.7, we validate this interpolation method.
### 2.5 Combining the clusters and the tidal field
We use the tidal tensor to calculate the external potential acting on each of
the stars in the simulated clusters. We integrate the internal potential of
the clusters with this external potential using a Bridge-like scheme (Fujii et
al., 2007), which is implemented in AMUSE (Portegies Zwart et al., 2013).
This scheme can be used to combine interacting systems that are calculated in
different instances and/or using different codes, i.e. multiple interacting
star clusters, globular clusters in a galactic environment, a galactic disk in
a halo potential or embedded star clusters (Pelupessy & Portegies Zwart,
2011).
In this scheme, the cluster experiences the gravity from the external field
through periodic velocity kicks. It alternates between these velocity kicks
and a drift due to self gravity evolution of the system. During one time step,
the system first experiences a kick of the velocities over a time step $dt/2$,
then a drift over a time step $dt$, and finally another kick over $dt/2$. In
our setup, the external potential is derived from the CosmoGrid simulation and
therefore necessarily fixed, while the clusters receive velocity kicks from
the external tidal field.
### 2.6 Escaping and bound stars
In the simulations, we calculate a tidal radius from the cluster mass and the
strength of the tidal field (given by the largest eigenvalue of the tidal
tensor $\lambda_{\rm max}$). This tidal radius is equal to
$R_{\rm tidal}=\left(\frac{GM}{\lambda_{\rm max}}\right)^{1/3}.$ (2)
Particles at a distance from the cluster centre larger than $R_{\rm tidal}$
will experience a larger force from the external tidal field than from the
cluster’s own internal mass. It is then considered an ‘escaping particle’, and
not included in the cluster’s bound mass (defined as the total mass inside
$R_{\rm tidal}$). If the particle returns to a position within the tidal
radius, this is reversed. If it moves to a distance $>10R_{\rm tidal}$, the
particle is removed from the simulation.
At times when the cluster is located in the centre of a local subconcentration
of dark matter or near a halocentre, the value of $\lambda_{\rm max}$ may
become negative for a short period. In such cases, the tidal radius is not
defined. At such times, the bound mass of the cluster is not evaluated and no
mass loss is experienced by the cluster. Stars that may have escaped from the
cluster during this period will however still be removed once the tidal radius
is again defined.
### 2.7 Validation
We validate the simulation environment by comparing our results with those
obtained using a different method and to a star cluster in isolation. As a
reference model, we simulate a star cluster in an orbit with $R_{\rm
apo}=12~{}h^{-1}{\rm kpc}$ and $e=0.71$ around a point-mass of $10^{10}$ ${\rm
M}_{\odot}$, by including this point-mass in the simulation. The simulated
cluster contains 8 000 equal-mass stars of 1 ${\rm M}_{\odot}$within a Plummer
(1911) sphere with a virial radius of 3 pc. We compare this model to a
simulation where we first calculated the tidal field that would be experienced
by such a cluster and used this as an external field for the simulation using
the Bridge scheme. For this test, we run simulations with ph4, as the large
difference in particle masses would make the reference simulation unsuitable
for a single-precision tree code like Bonsai.
The mass evolution of these test simulations are presented in Fig. 8. The
difference in the mass evolution of the cluster in the simulations with and
without Bridge is quite small, indicating the validity of this method for this
set of parameters.
Also, we validate the effect of the discretisation at which the tidal tensor
is evaluated. Since the time-resolution of the cosmological simulation is
limited to about 35 Myr at high $z$ and 17.5 Myr at low $z$, anomalies will be
visible for clusters with an orbital period of this order. Ideally, one would
like to increase the number of snapshots for the cosmological simulation in
order to obtain a higher time-resolution for the tidal tensor.
In Fig. 9 we present the mass evolution of a cluster in a static halo
potential (Paczynski, 1990) with a core mass of $10^{9}$ ${\rm M}_{\odot}$ and
a core radius of $1$ kpc, now using tidal tensors sampled with a time
resolution ranging from 1 to 35 Myr. The orbital parameters of the cluster in
this potential are $R_{\rm apo}=15$ kpc, $e=0.85$, and the cluster orbits the
potential in $670$ Myr. The mass loss rates for the clusters using a tidal
tensor sampled with a time interval of 1 Myr and 9 Myr have converged in this
orbit, while the clusters using tensors sampled with a 17.5 or 35 Myr time
interval show a slightly reduced mass loss rate and therefore longer
lifetimes. For clusters with closer orbits, this effect will be stronger. In a
cosmological setting, any change in the potential that takes place on a
timescale similar to or smaller than the sampling rate cannot be taken into
account properly, and may also lead to errors in the mass loss rate.
Figure 8: Mass evolution of a 8 000 star cluster in isolation (solid curve),
and in orbit around a point mass of $10^{10}$ ${\rm M}_{\odot}$(dashed and
dotted curves). The dashed curve gives the mass-evolution of the cluster when
the tidal field is incorporated in the gravitational $N$-body simulation using
Bridge. The tidal field in the latter case was resolved at the resolution of
the $N$-body integrator. The dotted curve gives the mass-evolution when the
point mass is directly included in the simulation. Figure 9: Mass evolution of
a 8 000 star cluster with tidal field calculated from a halo potential using
the tidal tensor and evaluated using Bridge at discrete time intervals. The
cluster orbits in the halo potential with $R_{\rm apo}=15\,{\rm kpc}$ and
$e=0.85$. In order to study the effect of the time resolution of the tidal
field, we evaluate the tidal tensor with time intervals of 1, 9, 18 and 35
Myr.
## 3 Results
### 3.1 The evolution of the selected haloes
We perform simulations of star clusters in two selected haloes, which we call
halo A and halo B. Halo A experiences two major merger events, the last of
which is not completed by $z=0$ (see Figs.1 and 3), but otherwise its
evolution is rather gradual from $z=65$ to $z=0$. By the end of the simulation
(at $z=0$) its mass is $6.3\times 10^{11}\,\mbox{$h^{-1}{\rm M}_{\odot}$}$.
Halo B has a more violent history; it experiences two major mergers, one
between 4 Gyr and 5 Gyr and one with comparable mass between 8 Gyr and 10 Gyr
(see also Figs. 1 and 3). By the end of the simulation at $z=0$ it has a mass
of $4.8\times 10^{11}$ $h^{-1}{\rm M}_{\odot}$ (see Tab. 1). Even though the
total mass of halo B is slightly smaller than that of halo A, both haloes
contain a mass of $7\times 10^{9}$ $h^{-1}{\rm M}_{\odot}$ within 3 kpc from
the halo centre (see Fig. 4).
In each halo we have selected 15 dark-matter particles which are initialized
at $z=10$ as star clusters, and evolved with the cosmological simulation as a
background potential.
### 3.2 The evolution of the star clusters
We perform 30 simulations of star clusters with a tidal field; 15 are
initialized in halo A and 15 in halo B. An additional cluster is simulated in
isolation, to identify the mass loss component caused by relaxation. For this
cluster, a radius of 200 pc is used to determine if stars are bound. All
clusters are born at $z=10$ (corresponding to an age of the Universe of about
500 Myr) with a total mass of $\sim 32\,000$ ${\rm M}_{\odot}$ and an initial
virial radius of 3 parsec for each cluster. All stars have the same mass and
we did not include stellar evolution.
In Fig. 10 and Tab. 2 we present the mass evolution and final masses of these
simulations. The mass evolution of all simulated clusters is rather gradual
irrespective of the sudden events in the growth of the host haloes. The small
and rather sudden changes in mass are caused by the pericentre passages of the
clusters in its orbit around the dark matter host. The clusters with a smaller
orbital separation at z=0 tend to lose mass at a higher rate.
The averaged final mass of the star clusters in halo A is smaller than the
final mass of those that evolved in halo B for each of the radial bins (see
Tab. 2). This is consistent with halo B originating from a larger number of
less massive haloes, causing the tidal forces experienced by clusters in this
halo over time to be smaller. For both haloes the clusters selected around
$3~{}h^{-1}{\rm kpc}$ show the strongest mass loss; this is noticeable from
the first few Gyr on. This behaviour is as expected from the distribution of
particles at $z=10$ (see Fig. 6), where we see that the particles that end up
in the central parts are already more concentrated at high redshift. The rate
of mass loss for these clusters proceeds more gradual compared to the clusters
in wider orbits. Integrated over time clusters lose mass at a rather constant
rate.
Figure 10: Bound mass in the simulated clusters for haloes A (top figure) and B (bottom figure). The red, solid; green, dashed and blue, dotted lines indicate clusters selected at $3$, $6$ and $12~{}h^{-1}{\rm kpc}$, respectively. The black line indicates the mass of a cluster without an external tidal field. Table 2: Results for star cluster simulations using Bonsai, with various tidal tensors. Each cluster contains 32 000 equal mass stars of about 1${\rm M}_{\odot}$, and has an initial radius of 3 parsec. We use a softening length of $40/N$. Cluster type is Immigrant (I) or Native (N). The clusters indicated with a bold font are displayed in Figs.12 and 14. Cluster ID | Halo | Type | Distance at z=0 | Mass at z=0
---|---|---|---|---
Isolated | | | | 28174 $\mathbf{{\rm\bf M}_{\odot}}$
A3-1 | A | I | $3~{}h^{-1}{\rm kpc}$ | 10157 ${\rm M}_{\odot}$
A3-2 | A | I | $3~{}h^{-1}{\rm kpc}$ | 15171 ${\rm M}_{\odot}$
A3-3 | A | N | $3~{}h^{-1}{\rm kpc}$ | 10969 ${\rm M}_{\odot}$
A3-4 | A | N | $3~{}h^{-1}{\rm kpc}$ | 11085 ${\rm M}_{\odot}$
A3-5 | A | N | $3~{}h^{-1}{\rm kpc}$ | 10169 ${\rm M}_{\odot}$
Average | | | | 11510 $\mathbf{{\rm\bf M}_{\odot}}$
A6-1 | A | I | $6~{}h^{-1}{\rm kpc}$ | 17906 ${\rm M}_{\odot}$
A6-2 | A | N | $6~{}h^{-1}{\rm kpc}$ | 24130 ${\rm M}_{\odot}$
A6-3 | A | N | $6~{}h^{-1}{\rm kpc}$ | 14964 ${\rm M}_{\odot}$
A6-4 | A | N | $6~{}h^{-1}{\rm kpc}$ | 17978 ${\rm M}_{\odot}$
A6-5 | A | N | $6~{}h^{-1}{\rm kpc}$ | 18430 ${\rm M}_{\odot}$
Average | | | | 18682 $\mathbf{{\rm\bf M}_{\odot}}$
A12-1 | A | I | $12~{}h^{-1}{\rm kpc}$ | 21677 ${\rm M}_{\odot}$
A12-2 | A | I | $12~{}h^{-1}{\rm kpc}$ | 16487 ${\rm M}_{\odot}$
A12-3 | A | N | $12~{}h^{-1}{\rm kpc}$ | 17574 ${\rm M}_{\odot}$
A12-4 | A | N | $12~{}h^{-1}{\rm kpc}$ | 15873 ${\rm M}_{\odot}$
A12-5 | A | N | $12~{}h^{-1}{\rm kpc}$ | 16904 ${\rm M}_{\odot}$
Average | | | | 17703 $\mathbf{{\rm\bf M}_{\odot}}$
B3-1 | B | I | $3~{}h^{-1}{\rm kpc}$ | 19422 ${\rm M}_{\odot}$
B3-2 | B | N | $3~{}h^{-1}{\rm kpc}$ | 18974 ${\rm M}_{\odot}$
B3-3 | B | N | $3~{}h^{-1}{\rm kpc}$ | 16366 ${\rm M}_{\odot}$
B3-4 | B | N | $3~{}h^{-1}{\rm kpc}$ | 15633 ${\rm M}_{\odot}$
B3-5 | B | N | $3~{}h^{-1}{\rm kpc}$ | 20839 ${\rm M}_{\odot}$
Average | | | | 18247 $\mathbf{{\rm\bf M}_{\odot}}$
B6-1 | B | I | $6~{}h^{-1}{\rm kpc}$ | 24774 ${\rm M}_{\odot}$
B6-2 | B | N | $6~{}h^{-1}{\rm kpc}$ | 18446 ${\rm M}_{\odot}$
B6-3 | B | I | $6~{}h^{-1}{\rm kpc}$ | 15988 ${\rm M}_{\odot}$
B6-4 | B | N | $6~{}h^{-1}{\rm kpc}$ | 21128 ${\rm M}_{\odot}$
B6-5 | B | N | $6~{}h^{-1}{\rm kpc}$ | 21398 ${\rm M}_{\odot}$
Average | | | | 20347 $\mathbf{{\rm\bf M}_{\odot}}$
B12-1 | B | I | $12~{}h^{-1}{\rm kpc}$ | 19986 ${\rm M}_{\odot}$
B12-2 | B | I | $12~{}h^{-1}{\rm kpc}$ | 23479 ${\rm M}_{\odot}$
B12-3 | B | N | $12~{}h^{-1}{\rm kpc}$ | 25819 ${\rm M}_{\odot}$
B12-4 | B | N | $12~{}h^{-1}{\rm kpc}$ | 26374 ${\rm M}_{\odot}$
B12-5 | B | I | $12~{}h^{-1}{\rm kpc}$ | 22137 ${\rm M}_{\odot}$
Average | | | | 23559 $\mathbf{{\rm\bf M}_{\odot}}$
In Tab. 2 we distinguish between two types of clusters: those that were part
of the main halo before the final completed major merger event (at $t=6$ Gyr
and $t=8$ Gyr for haloes A and B, respectively), and those that are accreted
upon or after this event. The former we identify with ‘native’ clusters, and
the others as ‘immigrants’. The difference between immigrant clusters and
native clusters is apparent in the Figs. 11.
(a) (b) (c) (d)
(e) (f)
Figure 11: Co-moving distance of the clusters to the main halo centre for
clusters in halo A (left) and halo B (right). Top, middle and bottom figures
show clusters selected at $3$, $6$ and $12~{}h^{-1}{\rm kpc}$, respectively.
Merger events in both haloes are visible as particles fall towards the halo
centre. The orbital periods of the clusters are clearly visible.
In Fig. 11 we present the orbital evolution of the selected dark-matter
particles (i.e. the clusters) from halo A (left) and halo B (right). In the
following paragraphs, we discuss the evolution of two clusters from each halo;
one immigrant and one native cluster in more detail.
#### 3.2.1 The clusters in halo A
In Fig. 12, we show the mass evolution, mass loss rate over intervals of 10
Myrs and tidal field strength for an immigrant cluster (nr. A6-1, left) and a
native cluster (nr. A6-4, right). This halo experiences a major merger which
starts at about 6 Gyr, at that time the immigrant cluster is also captured by
the main halo. The merger lasts until about 8 Gyr (see also Fig. 1). When the
merger is finished, the tidal field strength $\lambda_{\rm max}$ shows more
frequent peaks, indicating it has a shorter orbital period than before the
merger. However, there is little change in the mass-loss rate.
The native cluster of halo A experiences the same merger but was already
member of the major halo. Its orbit becomes somewhat less eccentric after the
merger, while its apocentric distance and orbital period decrease. The mass-
loss rate from this cluster is mostly unaffected by the merger.
(a) (b)
Figure 12: Evolution of two star clusters in orbit around halo A. To the left
is a typical immigrant cluster identified with dark-matter particle nr A6-1
and to the right we present a native cluster, particle nr A6-4 (see also
Tab.2). From top to bottom the panels show the bound mass, the mass-loss rate,
the strength of the tidal field ($\lambda_{\rm max}$) and the co-moving
distance of the cluster to the centre of the dark matter halo. The decrease
seen in the co-moving distance, apart from the merger event, is caused by the
expanding cosmic volume.
In Fig. 13, we show the result of two star clusters (A12-1 and A12-4), both
using the tidal tensor calculated from the evolving CosmoGrid halo and using a
tidal tensor calculated from the static $z=0$ CosmoGrid halo. For the static
halo case, we sampled the tensor using the orbital trajectory of the cluster
around its parent halo. For the native cluster, the resulting mass evolution
differs only marginally, while the immigrant cluster suffers considerably
higher mass loss in the static halo case, especially around the time the halo
merger takes place, when its orbit is erratic. The static halo is a good
approximation for the native cluster, while it falls short for the immigrant
cluster.
Figure 13: Evolution of two star clusters, A12-1 (immigrant) and A12-4
(native). Both are simulated in an evolving halo as well as a static halo.
#### 3.2.2 The clusters in halo B
In Fig. 14 we present the details of the evolution of an immigrant cluster
(nr. B6-3, left) and a native cluster (nr. B6-2, right) of halo B. The
immigrant cluster is captured during the major merger event that starts around
$t=8$ Gyr. Different from the clusters in halo A, the orbital period of the
cluster around the dark matter halo is hardly visible in the cluster’s mass
evolution (see Fig. 14, left, top panel). Just before the merger (around $t=7$
Gyr) the mass-loss rate is slightly smaller than before or after the merger.
This is caused by the distortion of the infalling halo of which this cluster
is a member at that time. When the merger is completed the cluster mass-loss
rate has resumed to be as high as before the merger (see Fig. 14, left, second
panel).
The native cluster becomes part of the main halo during its first major
merger, at around $t=5$ Gyr. After this merger, the tidal forces experienced
by the cluster are stronger than before, visible in Fig. 14 (right, third
panel) as a sudden increase of $\lambda_{\rm max}$ after $t=6$ Gyr. Its mass-
loss rate is also increased, as can be seen from Fig. 14 (right, second panel)
and the difference in slope of $M_{\rm bound}(t)$ at $t=4$ Gyr and $t=8$ Gyr.
The second major merger event leads to a temporary reduction in mass-loss for
the cluster (at around $t=10$ Gyr). After this second halo merger the orbital
period is evidently visible in the mass evolution of the cluster.
(a) (b)
Figure 14: Evolution of two star clusters in orbit around halo B. To the left
is a typical immigrant cluster identified with dark-matter particle nr B6-3
and to the right we present a native cluster, particle nr B6-2 (see also
Tab.2). From top to bottom the panels show the bound mass, the mass-loss rate,
the strength of the tidal field ($\lambda_{\rm max}$) and the co-moving
distance of the cluster to the centre of the dark matter halo. Negative values
of $\lambda_{\rm max}$ correspond to times where the cluster is located inside
a local (sub)concentration of matter. The spikes in the bottom panels for the
immigrant and the native clusters is associated with a confusion in
identifying the main parent in the halo finder.
## 4 Discussion and conclusions
We presented a method to simulate star clusters within a pre-calculated tidal
field, using the AMUSE environment. As a proof of concept, we apply this
method to calculate the mass loss rates for star clusters in two live
$\Lambda$CDM haloes from the CosmoGrid simulation. Our method compares well to
self-consistent simulations.
We find that the mass loss rate strongly depends on the cluster’s orbital
parameters around the halo centre, as well as the central mass of the halo.
Also, tidal disruption due to the $\Lambda$CDM environment is weaker in haloes
that experienced many mergers. Finally, we find that in a Milky Way-like halo,
the contribution of the $\Lambda$CDM environment to $\Delta M/M_{\rm init}$
can be up to 0.6 for clusters with an initial mass of 32 000${\rm M}_{\odot}$.
In this article, several effects have not been taken into account, most
notably we used a dark matter-only simulation. In this section, we discuss the
relevance of these effects and how we intend to address these shortcomings.
Since the CosmoGrid simulation is a dark matter-only simulation, we do not
account for the effect of baryons. Because the formation of the large-scale
environment is dominated by dark matter, the lack of baryons has little
influence on the formation of the haloes. Star clusters however likely contain
little or no dark matter (Baumgardt et al., 2010; Conroy, Loeb & Spergel,
2011).
Because our simulation lacks baryons, there is no indication where star
clusters would have formed or ended up had they been formed in our
cosmological environment. We therefore have to resort to our method of
identifying tracer particles for the star clusters in the final snapshot. We
selected particles based on their distance to the halo centre. However, it is
very likely that globular clusters would have formed in more specific
locations, and possibly followed paths quite different from the ones in our
simulation. In a simulation that includes baryonic matter and star formation,
it would be possible to detect locations and masses of star clusters as they
form.
The difference between old and young clusters herein is large: old clusters
(like the globular clusters in the Milky Way) formed before there was a
Galactic disk, and remain relatively free of its influence. Young clusters
however form in the galactic disk, and the tides experienced by these clusters
are dominated by encounters with giant molecular clouds and spiral arms
(Lamers & Gieles, 2006; Gieles, Athanassoula & Portegies Zwart, 2007), the
effect of which is about four times larger than the tidal field (e.g. Lamers &
Gieles, 2006, Figure 1). For this reason, we focus on old stellar clusters and
initialize our simulated star clusters at an early epoch, before the galactic
environment would have formed. In order to simulate young clusters, a galaxy
simulation including baryons would be required. However, when the old clusters
formed, the GMC density in the star-forming environment was likely very high,
causing early disruption of low-mass globular clusters (Elmegreen, 2010). This
effect is not included in our simulations.
In order to investigate a large number of clusters with a reasonable amount of
stars, we used the Barnes & Hut (1986) tree code Bonsai for most of our star
cluster simulations. We compared the results for two distinct star clusters to
similar simulations with the direct $N$-body code ph4. The results for both
runs are similar for both codes, the direct code showing enhanced mass loss
around the time of core collapse. However, the tree code fails to accurately
describe the inner structure of the star cluster, and requires the distance
between stars to be softened.
Another limit of our environment is its resolution. The spatial resolution of
CosmoGrid is given by the softening length employed (175 parsec). Forces that
occur on a scale similar to or smaller than this softening length are not
accurately taken into account. Likewise, the mass resolution of dark matter
particles in CosmoGrid is $1.28\times 10^{5}$ ${\rm M}_{\odot}$, about 4 times
larger than the initial mass of our simulated clusters. The effects of tidal
forces caused by a small dark matter object passing at close range to our
clusters (such as a subhalo) are therefore limited, creating a possible bias
against the effect of such structures. However, this resolution effect would
be more important for baryons than it is for dark matter.
The benefit of using a large-scale cosmological simulation however, is that
the formation of a dark-matter halo is followed. A simulation that only models
a collision between galaxies would not take the earlier history and
distribution of star clusters into account. In order to have the benefits of
both a cosmological environment and high resolution, one could use re-
simulation, where galaxies are simulated at high resolution within a lower-
resolution environment.
Another limit imposed by our use of a pre-calculated simulation is its limited
number of snapshots. We lack continuous information about the tidal tensor. In
order to prevent sudden changes in the tidal field, we interpolate the tidal
tensor between snapshots. However, it remains impossible to accurately track
sudden changes in the tidal field on timescales shorter than our time
resolution, such as those occurring during halo mergers. Since the orbital
periods of our clusters around the halo centre are larger than this time
resolution, we do not expect this to have a large influence on the evolution
of the tidal field. However, short-lasting passages of nearby objects may not
be taken into account accurately, and a method in which the tidal field is
sampled at more intervals remains preferable.
In a follow-up article (Rieder et al., in prep.), we will apply the method
described in this article to the evolution of star clusters in the disk of a
simulated Milky Way-type galaxy. In this follow-up, we will address several of
the limitations discussed above, especially the lack of baryons and as a
result the orbits and origins of the clusters.
## Acknowledgements
It is our pleasure to thank the anonymous referee for very helpful suggestions
and comments that greatly helped to improve the article. Also, we are grateful
to Arjen van Elteren, Derek Groen, Inti Pelupessy, Jeroen Bédorf, Mark Gieles
and Nathan de Vries for support, interesting discussions and useful
suggestions.
This work was supported by NWO (grants IsFast [#643.000.803], VICI
[#639.073.803], LGM [#612.071.503] and AMUSE [#614.061.608]), NCF (grants
[#SH-095-08] and [#SH-187-10]), NOVA and the LKBF in the Netherlands, and by
NSF grant AST-0708299 in the U.S. T.I. is financially supported by MEXT HPCI
STRATEGIC PROGRAM and MEXT/JSPS KAKENHI Grant Number 24740115. We thank the
DEISA Consortium (EU FP6 project RI-031513 and FP7 project RI-222919) for
support within the DEISA Extreme Computing Initiative (GBBP project).
The Cosmogrid simulations were partially carried out on Cray XT4 at Center for
Computational Astrophysics, CfCA, of National Astronomical Observatory of
Japan; Huygens at the Dutch National High Performance Computing and e-Science
Support Center, SARA (The Netherlands); HECToR at the Edinburgh Parallel
Computing Center (United Kingdom) and Louhi at IT Center for Science in Espoo
(Finland).
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|
arxiv-papers
| 2013-09-27T20:00:15 |
2024-09-04T02:49:51.612070
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Steven Rieder, Tomoaki Ishiyama, Paul Langelaan, Junichiro Makino,\n Stephen L. W. McMillan and Simon Portegies Zwart",
"submitter": "Steven Rieder",
"url": "https://arxiv.org/abs/1309.7352"
}
|
1309.7446
|
# estimates of the gaps between consecutive eigenvalues of Laplacian
Daguang Chen∗, Tao Zheng, Hongcang Yang∗∗
###### Abstract.
By the calculation of the gap of the consecutive eigenvalues of
$\mathbb{S}^{n}$ with standard metric, using the Weyl’s asymptotic formula, we
know the order of the upper bound of this gap is $k^{\frac{1}{n}}.$ We
conjecture that this order is also right for general Dirichlet problem of the
Laplace operator, which is optimal if this conjecture holds, obviously. In
this paper, using new method, we solve this conjecture in the Euclidean space
case intrinsically. We think our method is valid for the case of general
Riemannian manifolds and give some examples directly.
###### Key words and phrases:
Laplacian, consecutive eigenvalues, test function, Riemannian manifold,
Hyperbolic space
###### 2010 Mathematics Subject Classification:
Primary 35P15, 58C40; Secondary 58J50.
*The work of the first named author was partially supported by NSFC grant No. 11101234.
**The work of the third named author was partially supported by NSFC and SF of
CAS.
## 1\. introduction
Let $\Omega$ be a bounded domain in an $n$-dimensional complete Riemannian
manifold $M$ with boundary (possible empty). Then the Dirichlet eigenvalue
problem of Laplacian on $\Omega$ is given by
$\left\\{\begin{aligned} \Delta u=&-\lambda u,\;\mbox{in}\;\Omega,\\\
u=&0,\;\;\;\;\;\;\;\;\mbox{on}\;\partial\Omega,\end{aligned}\right.$ (1.1)
where $\Delta$ is Laplaican on $M$. It is well known that the spectrum of
(1.1) has the real and purely discrete eigenvalues
$0<\lambda_{1}<\lambda_{2}\leq\lambda_{3}\leq\cdots\nearrow\infty,$ (1.2)
where each $\lambda_{i}$ has finite multiplicity which is repeated according
to its multiplicity. The corresponding orthonormal basis of real
eigenfunctions will be denoted $\\{u_{j}\\}_{j=1}^{\infty}$. We go forward
under the assumption that $L^{2}(\Omega)$ represents the real Hilbert space of
real-valued $L^{2}$ functions on $\Omega$. We put $\lambda_{0}=0$ if
$\partial\Omega=\emptyset$.
An important aspect of estimating higher eigenvalues is to obtain as precise
as possible the estimate of gaps of consecutive eigenvalues of (1.1). In this
regard, we will review some important results on the estimates of eigenvalue
problem (1.1).
For the upper bound of the gap of consecutive eigenvalues of (1.1), when
$\Omega$ is a bounded domain in an $2$-dimensional Euclidean space
$\mathbb{R}^{2},$ in 1956, Payne, Pólya and Weinberger (cf.[29] and [30])
proved
$\lambda_{k+1}-\lambda_{k}\leq\frac{2}{k}\sum_{i=1}^{k}\lambda_{i}.$ (1.3)
C. J. Thompson [34], in 1969, extended (1.3) to $n$-dimensional case and
obtained
$\lambda_{k+1}-\lambda_{k}\leq\frac{4}{nk}\sum_{i=1}^{k}\lambda_{i}.$ (1.4)
Hile and Protter [25] improved (1.4) to
$\sum_{i=1}^{k}\frac{\lambda_{i}}{\lambda_{k+1}-\lambda_{i}}\geq\frac{nk}{4}.$
(1.5)
Yang (cf. [37] and more recently [18]) has obtained a sharp inequality
$\sum_{i=1}^{k}(\lambda_{k+1}-\lambda_{i})\left(\lambda_{k+1}-\left(1+\frac{4}{n}\right)\lambda_{i}\right)\leq
0.$ (1.6)
From (1.6), one can infer
$\lambda_{k+1}\leq\frac{1}{k}\left(1+\frac{4}{n}\right)\sum_{i=1}^{k}\lambda_{i}.$
(1.7)
The inequalities (1.6) and (1.7) are called Yang’s first inequality and second
inequality, respectively (cf. [8, 9, 7, 23]). Also we note that Ashbaugh and
Benguria gave an optimal estimate for $k=1$ (cf. [2, 3, 4]). From the
Chebyshev’s inequality, it is easy to prove the following relations
$(\ref{CY41})\Longrightarrow(\ref{CY42})\Longrightarrow(\ref{hp1})\Longrightarrow(\ref{ppwtype}).$
From (1.6), Cheng and Yang [16] obtained
$\lambda_{k+1}-\lambda_{k}\leq
2\left[\left(\frac{2}{n}\frac{1}{k}\sum\limits_{i=1}^{k}\lambda_{i}\right)^{2}-\left(1+\frac{4}{n}\right)\frac{1}{k}\sum\limits_{i=1}^{k}\left(\lambda_{i}-\frac{1}{k}\sum\limits_{j=1}^{k}\lambda_{j}\right)^{2}\right]^{\frac{1}{2}}.$
(1.8)
Cheng and Yang [18], using their recursive formula, obtained
$\lambda_{k+1}\leq C_{0}(n)k^{\frac{2}{n}}\lambda_{1},$ (1.9)
where $C_{0}(n)\leq 1+\frac{4}{n}$ is a constant (see Cheng and Yang’s paper
[18]). From the Weyl’s asymptotic formula (cf. [35]), we know that the upper
bound (1.9) of Cheng and Yang is best possible in the meaning of the order on
$k$.
For a complete Riemannian manifold $M$, from the Nash’s theorem [26], there
exists an isometric immersion
$\psi\,:\;M\longrightarrow\mathbb{R}^{N},$
where $\mathbb{R}^{N}$ is Euclidean space. The mean curvature of the immersion
$\psi$ is denoted by $H$ and $|H|$ denotes its norm. Define
$\Phi=\\{\psi\;|\psi\;\mbox{is an isometric immersion from $M$ into Euclidean
space}\\}.$
When $\Omega$ is a bounded domain of a complete Riemannian manifold $M$,
isometrically immersed into a Euclidean space $\mathbb{R}^{N}$, Cheng and the
first author [14] (cf. [20, 24])obtained
$\sum_{i=1}^{k}(\lambda_{k+1}-\lambda_{i})^{2}\leq\frac{4}{n}\sum_{i=1}^{k}(\lambda_{k+1}-\lambda_{i})\left(\lambda_{i}+\frac{n^{2}}{4}H_{0}^{2}\right),$
(1.10)
where
$H_{0}^{2}=\inf_{\psi\in\Phi}\sup_{\Omega}|H|^{2}.$ (1.11)
Using the recursive formula in Cheng and Yang [18], Cheng and the first author
in [14] also deduced
$\lambda_{k+1}+\frac{n^{2}}{4}H_{0}^{2}\leq
C_{0}(n)k^{\frac{2}{n}}\left(\lambda_{1}+\frac{n^{2}}{4}H_{0}^{2}\right),$
(1.12)
where $H_{0}^{2},C_{0}(n)$ are given by (1.11) and (1.9) respectively.
From (1.10), we can get the gaps of the consecutive eigenvalues of Laplacian
$\lambda_{k+1}-\lambda_{k}\leq
2\left(\left(\frac{2}{n}\frac{1}{k}\sum\limits_{i=1}^{k}\lambda_{i}+\frac{n}{2}H_{0}^{2}\right)^{2}-\left(1+\frac{4}{n}\right)\frac{1}{k}\sum\limits_{i=1}^{k}\left(\lambda_{i}-\frac{1}{k}\sum\limits_{j=1}^{k}\lambda_{j}\right)^{2}\right)^{\frac{1}{2}}.$
(1.13)
###### Remark 1.1.
When $\Omega$ is an $n$-dimensional compact homogeneous Riemannian manifold, a
compact minimal submanifold without boundary and a connected bounded domain in
the standard unit sphere $\mathbb{S}^{N}(1)$, and a connected bounded domain
and a compact complex hypersurface without boundary of the complex projective
space $\mathbb{CP}^{n}(4)$ with holomorphic sectional curvature 4, many
mathematicians have sutudied the universal inequalities for eigenvalues and
the difference of the consecutive eigenvalues (cf. [16, 17, 19, 21, 22, 23,
28, 38, 27, 33]).
###### Remark 1.2.
Another problem is the lower bound of the gap of the first two eigenvalues. In
general, there exists the famous fundamental gap conjecture for the Dirichlet
eigenvalue problem of the Schrödinger operator(cf.[11, 1, 39, 32, 40]and the
references therein). The fundamental gap conjecture was solved by B. Andrews
and J. Clutterbuck in [10] .
From (1.8) and (1.13), it is not difficult to see that both Yang’s estimate
for the gap of consecutive eigenvalues of (1.1) implicited in [37] and the
estimate from [14] are on the order of $k^{\frac{3}{2n}}$. However, by the
calculation of the gap of the consecutive eigenvalues of $\mathbb{S}^{n}$ with
standard metric, using the Weyl’s asymptotic formula, we know the order of the
upper bound of this gap is $k^{\frac{1}{n}}.$ Hence we conjecture that
###### Conjecture 1.1.
Let $\Omega$ be a bounded domain in an $n$-dimensional complete Riemannian
manifold $M$.For the Dirichlet problem (1.1), the upper bound for the gap of
consecutive eigenvalues of Laplacian should be
$\lambda_{k+1}-\lambda_{k}\leq C_{n,\Omega}k^{\frac{1}{n}},$ (1.14)
where $C_{n,\Omega}$ is a constant dependent on $\Omega$ itself and the
dimension $n$.
###### Remark 1.3.
The famous Panye-Pólya-Weinberger conjecture(cf.[29, 30, 34, 5, 6])claims
that, when $M=\mathbb{R}^{n}$, for Dirichlet eigenvalue problem (1.1), one
should have
$\dfrac{\lambda_{k+1}}{\lambda_{k}}\leq\left.\dfrac{\lambda_{2}}{\lambda_{1}}\right|_{\mathbb{B}^{n}}=\left(\dfrac{j_{n/2,1}}{j_{n/2-1,1}}\right)^{2},$
(1.15)
where $\mathbb{B}^{n}$ is the $n$-dimensional unit ball in $\mathbb{R}^{n},$
and $j_{p,k}$ is the $k^{th}$ positive zero of the Bessel function $J_{p}(t)$.
From the Weyl’s asymptotic formula and (1.15), we know that the order of the
upper bounder of the consective eigenvalues of eigenvalue problem (1.1) is
just $k^{\frac{2}{n}}$. This is why we make this conjecture.
In the following, the constant $C_{n,\Omega}$ are allowed to be different in
different cases.
When $\Omega$ is a bouded domain in $\mathbb{R}^{n}$, for the Dirichlet
eigenvalue problem (1.1), we give the affirmative answer to the conjecture
(1.14).
###### Theorem 1.1.
Let $\Omega\subset\mathbb{R}^{n}$ be a bounded domain in Euclidean space
$\mathbb{R}^{n}$ and $\lambda_{k}$ be the $k^{th}$ eigenvalue of the Dirichlet
eigenvalue problem (1.1). Then we have
$\lambda_{k+1}-\lambda_{k}\leq C_{n,\Omega}k^{\frac{1}{n}},$ (1.16)
where $C_{n,\Omega}=4\lambda_{1}\sqrt{\frac{C_{0}(n)}{n}}$, $C_{0}(n)$ is
given by $(\ref{diffyong})$.
We think our method of proving Theorem 1.1 will be valid for the case of
$n$-dimensional complete Riemannian manifold. Here, we give some examples
directly.
###### Corollary 1.2.
Let $\Omega\subset\mathbb{H}^{n}(-1)$ be a bounded domain in hyperbolic space
$\mathbb{H}^{n}(-1)$, and $\lambda_{k}$ be the $k^{th}$ eigenvalue of the
Dirichlet eigenvalue problem (1.1). Then we have
$\lambda_{k+1}-\lambda_{k}\leq C_{n,\Omega}k^{\frac{1}{n}},$ (1.17)
where $C_{n,\Omega}$ depends on $\Omega$ and the dimension $n$, given by
$C_{n,\Omega}=4\left[C_{0}(n)\left(\lambda_{1}-\frac{(n-1)^{2}}{4}\right)\Big{(}\lambda_{1}+\frac{n^{2}}{4}H_{0}^{2}\Big{)}\right]^{\frac{1}{2}},$
(1.18)
$C_{0}(n)$ and $H_{0}^{2}$ are the same as the ones in (1.12).
In fact, by the comparison theorem for the distance function in Riemannian
manifold, we have
###### Corollary 1.3.
Let $\Omega\subset M$ be a bounded domain of an $n$-dimensional ($n\geq 3$)
simply connected complete noncompact Riemannian manifold $M$ with sectional
curvature $Sec$ satisfying
$-a^{2}\leq Sec\leq-b^{2},$
where $a\geq b\geq 0$ are constants. Let $\lambda_{k}$ be the $k^{th}$
eigenvalue of the eigenvalue problem (1.1). Then we have
$\lambda_{k+1}-\lambda_{k}\leq C_{n,\Omega}k^{\frac{1}{n}}$ (1.19)
where $C_{n,\Omega}$ depends on $\Omega$ and the dimension $n$, given by
$C_{n,\Omega}=4\left[C_{0}(n)\left(\lambda_{1}-\frac{(n-1)^{2}}{4}b^{2}+\frac{a^{2}-b^{2}}{4}\right)\left(\lambda_{1}+\frac{n^{2}}{4}H_{0}^{2}\right)\right]^{\frac{1}{2}},$
(1.20)
$C_{0}(n)$ and $H_{0}^{2}$ are the same as the ones in (1.12).
###### Remark 1.4.
Under the same assumption of Corollary 1.3, Lu and the first two authors [15]
obtained
$\sum_{i=1}^{k}(\lambda_{k+1}-\lambda_{i})^{2}\leq
4\sum_{i=1}^{k}(\lambda_{k+1}-\lambda_{i})\left(\lambda_{i}-\frac{(n-1)^{2}}{4}b^{2}+\frac{n-1}{2}(a^{2}-b^{2})\right).$
Therefore, one can get
$\displaystyle\lambda_{k+1}-\lambda_{k}\leq 2$
$\displaystyle\left[\left(\frac{2}{k}\sum_{i=1}^{k}\lambda_{i}-\frac{(n-1)^{2}}{4}b^{2}+\frac{n-1}{2}(a^{2}-b^{2})\right)^{2}\right.$
$\displaystyle\left.\quad-\frac{5}{k}\sum_{i=1}^{k}\left(\lambda_{i}-\frac{1}{k}\sum_{i=1}^{k}\lambda_{j}\right)^{2}\right]^{\frac{1}{2}}.$
## 2\. proofs of main results
In this section, we will give the proof of Theorem 1.1. In order to prove our
main results, we need the following key lemma
###### Lemma 2.1.
For the Dirichlet eigenvalue problem (1.1), let $u_{k}$ be the orthonormal
eigenfunction corresponding to the $k^{th}$ eigenvalue $\lambda_{k}$, i.e.
$\left\\{\begin{aligned} \Delta
u_{k}=&-\lambda_{k}u_{k},\;\text{in}\;\Omega,\\\
u_{k}=&0,\;\;\;\;\;\;\;\;\;\;\;\text{on}\;\partial\Omega,\\\
\int_{\Omega}u_{i}u_{j}&=\delta_{ij}.\end{aligned}\right.$
Then for any complex value function $g\in C^{3}(\Omega)\cap
C^{2}(\overline{\Omega})$ and $k,\,i\in\mathbb{Z}^{+},\,k>i\geq 1,$ we have
$\displaystyle\Big{(}(\lambda_{k+1}-\lambda_{i})+(\lambda_{k+2}-\lambda_{i})\Big{)}\int_{\Omega}|\nabla
g|^{2}u_{i}^{2}\leq$ $\displaystyle\int_{\Omega}\Big{|}2\nabla g\cdot\nabla
u_{i}+u_{i}\Delta g\Big{|}^{2}$ (2.1)
$\displaystyle+(\lambda_{k+1}-\lambda_{i})(\lambda_{k+2}-\lambda_{i})\int_{\Omega}|gu_{i}|^{2}.$
###### Proof.
For $i<k$, define
$\left\\{\begin{aligned} a_{ij}=&\int_{\Omega}gu_{i}u_{j},\\\
b_{ij}=&\int_{\Omega}\left(\nabla g\cdot\nabla u_{i}+\frac{1}{2}u_{i}\Delta
g\right)u_{j},\\\
\varphi_{i}=&gu_{i}-\sum_{j=1}^{k}a_{ij}u_{j},\end{aligned}\right.$
where $\nabla$ denotes the gradient operator. Obviously,
$a_{ij}=a_{ji},\;\;\int_{\Omega}\varphi_{i}u_{j}=0,\;\;\mbox{for}\;\;j=1,2,\cdots,k.$
(2.2)
Then, from the Stokes’ theorem, we get
$\lambda_{j}a_{ij}=\int_{\Omega}gu_{i}(-\Delta
u_{j})=-\int_{\Omega}(u_{i}\Delta g+g\Delta u_{i}+2\nabla g\cdot\nabla
u_{i})u_{j},$
i.e.
$2b_{ij}=(\lambda_{i}-\lambda_{j})a_{ij}.$ (2.3)
From the Stokes’ theorem, we have
$-2\int_{\Omega}gu_{i}\nabla\overline{g}\cdot\nabla
u_{i}=-\int_{\Omega}g\nabla\overline{g}\cdot\nabla
u_{i}^{2}=\int_{\Omega}(\nabla
g\cdot\nabla\overline{g}+g\Delta\overline{g})u_{i}^{2}.$
By the definition of $a_{ij}$, $b_{ij}$ and (2.3), we obtain
$\int_{\Omega}|\nabla
g|^{2}u_{i}^{2}=-2\int_{\Omega}gu_{i}\left(\nabla\overline{g}\cdot\nabla
u_{i}+\frac{1}{2}u_{i}\Delta\overline{g}\right)=-2\sum_{j=1}^{\infty}a_{ij}\overline{b_{ij}}=\sum_{j=1}^{\infty}(\lambda_{j}-\lambda_{i})|a_{ij}|^{2}.$
Similarly, from the Stokes’ theorem, (2.2) and (2.3), we have
$\displaystyle\int_{\Omega}|\nabla\varphi_{i}|^{2}=$
$\displaystyle\int_{\Omega}\overline{\varphi_{i}}\left(-\Delta\varphi_{i}\right)$
(2.4) $\displaystyle=$
$\displaystyle\int_{\Omega}\overline{\varphi_{i}}\left(-\Delta(gu_{i})+\sum\limits_{j=1}^{k}a_{ij}\Delta
u_{j}\right)$ $\displaystyle=$
$\displaystyle\int_{\Omega}\overline{\varphi_{i}}\left(-2\nabla g\cdot\nabla
u_{i}-\Delta gu_{i}-g\Delta u_{i}\right)$ $\displaystyle=$
$\displaystyle-\int_{\Omega}\overline{\varphi_{i}}\left(2\nabla g\cdot\nabla
u_{i}+\Delta gu_{i}-\lambda_{i}gu_{i}\right)$ $\displaystyle=$
$\displaystyle-2\sum\limits_{j=k+1}^{\infty}\overline{a_{ij}}b_{ij}+\lambda_{i}\sum\limits_{j=k+1}^{\infty}|a_{ij}|^{2}$
$\displaystyle=$
$\displaystyle\sum\limits_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})|a_{ij}|^{2}+\lambda_{i}\sum\limits_{j=k+1}^{\infty}|a_{ij}|^{2}.$
From the Rayleigh-Ritz inequality (cf. [13]) and (2.4), we have
$\lambda_{k+1}\leq\frac{\int_{\Omega}|\nabla\varphi_{i}|^{2}}{\int_{\Omega}|\varphi_{i}|^{2}}=\frac{\sum_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})|a_{ij}|^{2}}{\sum_{j=k+1}^{\infty}|a_{ij}|^{2}}+\lambda_{i},$
i.e.
$(\lambda_{k+1}-\lambda_{i})\sum_{j=k+1}^{\infty}|a_{ij}|^{2}\leq\sum_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})|a_{ij}|^{2}.$
(2.5)
From the Cauchy-Schwarz inequality, we have
$\left(\sum_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})|a_{ij}|^{2}\right)^{2}\leq\sum_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})^{2}|a_{ij}|^{2}\sum_{j=k+1}^{\infty}|a_{ij}|^{2},$
i.e.
$\displaystyle\left(\int_{\Omega}|\nabla
g|^{2}u_{i}^{2}-\sum_{j=1}^{k}(\lambda_{j}-\lambda_{i})|a_{ij}|^{2}\right)^{2}$
(2.6) $\displaystyle\leq$
$\displaystyle\left(\int_{\Omega}|gu_{i}|^{2}-\sum_{j=1}^{k}|a_{ij}|^{2}\right)\left(\int_{\Omega}\left|2\nabla
g\cdot\nabla u_{i}+u_{i}\Delta
g\right|^{2}-\sum_{j=1}^{k}(\lambda_{j}-\lambda_{i})^{2}|a_{ij}|^{2}\right).$
Define
$\left\\{\begin{aligned}
\widetilde{B}(i)=&\int_{\Omega}|gu_{i}|^{2}-\sum_{j=1}^{k}|a_{ij}|^{2}=\sum_{j=k+1}^{\infty}|a_{ij}|^{2}>0,\\\
\widetilde{A}(i)=&\int_{\Omega}\left|2\nabla g\cdot\nabla u_{i}+u_{i}\Delta
g\right|^{2}-\sum_{j=1}^{k}(\lambda_{j}-\lambda_{i})^{2}|a_{ij}|^{2}=\sum_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})^{2}|a_{ij}|^{2}\geq
0,\\\ \widetilde{C}(i)=&\int_{\Omega}|\nabla
g|^{2}u_{i}^{2}-\sum_{j=1}^{k}(\lambda_{j}-\lambda_{i})|a_{ij}|^{2}=\sum_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})|a_{ij}|^{2}.\end{aligned}\right.$
Next, we will deduce the maxima of $\widetilde{C}(i)$ by using the Lagrange
method of multipliers (cf.[12]). For any sequence
$\\{\phi_{ij}\\}_{j=k+1}^{\infty}$ satisfying
$\sum_{j=k+1}^{\infty}|\phi_{ij}|^{2}\leq\infty$
we define the function,
$\displaystyle\Phi(|\phi_{ij}|,\mu,\lambda)=$
$\displaystyle\sum_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})|\phi_{ij}|^{2}+\mu\left(\sum_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})^{2}|\phi_{ij}|^{2}-\widetilde{A}(i)\right)$
$\displaystyle+\lambda\left(\sum_{j=k+1}^{\infty}|\phi_{ij}|^{2}-\widetilde{B}(i)\right).$
where $\mu,\lambda$ are two real parameters.
Assume $\\{a_{ij}\\}_{j=k+1}^{\infty}$ is the extreme point of $\varphi$. Then
for any $\\{\psi_{ij}\\}_{j=k+1}^{\infty}$ satisfying
$\sum\limits_{j=k+1}^{\infty}|\psi_{ij}|^{2}<\infty$, from
$\left.\frac{\mathrm{d}}{\mathrm{d}t}\right|_{t=0}\Phi(|a_{ij}|+t|\psi_{ij}|)=0$
we have
$2\sum\limits_{j=k+1}^{\infty}|a_{ij}||\psi_{ij}|\left((\lambda_{j}-\lambda_{i})+\mu(\lambda_{j}-\lambda_{i})^{2}+\lambda\right)=0.$
(2.7)
Taking
$\psi_{ij}=\left\\{\begin{aligned} 1,&\;\;j=p,\\\
0,&\;\;\mbox{otherwise},\end{aligned}\right.$
in (2.7), we have
$|a_{ip}|\left((\lambda_{p}-\lambda_{i})+\mu(\lambda_{p}-\lambda_{i})^{2}+\lambda\right)=0,\;p=k+1,\cdots.$
(2.8)
From
$\left\\{\begin{aligned} &\frac{\partial\Phi}{\partial\mu}=0,\\\
&\frac{\partial\Phi}{\partial\lambda}=0,\end{aligned}\right.$
we have the two constraint conditions
$\left\\{\begin{aligned}
&\sum_{j=k+1}^{\infty}(\lambda_{j}-\lambda_{i})^{2}|a_{ij}|^{2}=\widetilde{A}(i),\\\
&\sum_{j=k+1}^{\infty}|a_{ij}|^{2}=\widetilde{B}(i).\end{aligned}\right.$
(2.9)
Since there are two Lagrange multipliers and $\widetilde{B}(i)>0,$ from (2.8)
and (2.9), there exist $r>l>k$ such that $|a_{ir}|\cdot|a_{il}|\neq
0,\;\lambda_{r}>\lambda_{l},$ and $|a_{ij}|=0,\;j\neq r,l.$ Hence, we have
$\left\\{\begin{aligned}
&m_{r}(\lambda_{r}-\lambda_{i})^{2}|a_{ir}|^{2}+m_{l}(\lambda_{l}-\lambda_{i})^{2}|a_{il}|^{2}=\widetilde{A}(i),\\\
&m_{r}|a_{ir}|^{2}+m_{l}|a_{il}|^{2}=\widetilde{B}(i),\end{aligned}\right.$
(2.10)
where $m_{r},m_{l}$ are the multiplicity of the eigenvalues
$\lambda_{r}\,\mbox{and}\,\lambda_{l},$ respectively. From (2.10), we have
$\widetilde{C}(i)=\frac{\widetilde{A}(i)+(\lambda_{r}-\lambda_{i})(\lambda_{l}-\lambda_{i})\widetilde{B}(i)}{(\lambda_{l}-\lambda_{i})+(\lambda_{r}-\lambda_{i})}.$
(2.11)
From (2.6), we have
$\widetilde{C}(i)\leq\sqrt{\widetilde{A}(i)\widetilde{B}(i)}.$ (2.12)
By the definition of $\widetilde{A}(i)$ and $\widetilde{B}(i),$ we have
$(\lambda_{k+1}-\lambda_{i})\leq\sqrt{\widetilde{A}(i)/\widetilde{B}(i)}.$
(2.13)
From the range of the function
$\frac{\widetilde{A}(i)+(\lambda_{r}-\lambda_{i})(\lambda_{l}-\lambda_{i})\widetilde{B}(i)}{(\lambda_{l}-\lambda_{i})+(\lambda_{r}-\lambda_{i})},$
(2.12), we have $r=k+2.$ From (2.13), we have $l=k+1.$ Therefore, we obtain
$\widetilde{C}(i)\leq\frac{\widetilde{A}(i)+(\lambda_{k+2}-\lambda_{i})(\lambda_{k+1}-\lambda_{i})\widetilde{B}(i)}{(\lambda_{k+2}-\lambda_{i})+(\lambda_{k+1}-\lambda_{i})}.$
(2.14)
From (2.14), and the definition of $\widetilde{A}(i),\;\widetilde{B}(i)$ and
$\widetilde{C}(i),$ we have
$\displaystyle\left((\lambda_{k+2}-\lambda_{i})+(\lambda_{k+1}-\lambda_{i})\right)\int_{\Omega}|\nabla
g|^{2}u_{i}^{2}$ (2.15) $\displaystyle\leq$
$\displaystyle\int_{\Omega}\left|2\nabla g\cdot\nabla u_{i}+u_{i}\Delta
g\right|^{2}+(\lambda_{k+1}-\lambda_{i})(\lambda_{k+2}-\lambda_{i})\int_{\Omega}|gu_{i}|^{2}$
$\displaystyle-\sum_{j=1}^{k}(\lambda_{k+1}-\lambda_{j})(\lambda_{k+2}-\lambda_{j})|a_{ij}|^{2}$
$\displaystyle\leq$ $\displaystyle\int_{\Omega}\left|2\nabla g\cdot\nabla
u_{i}+u_{i}\Delta
g\right|^{2}+(\lambda_{k+1}-\lambda_{i})(\lambda_{k+2}-\lambda_{i})\int_{\Omega}|gu_{i}|^{2}$
which finishes the proof of Lemma 2.1. ∎
Based on Lemma 2.1, we have
###### Corollary 2.2.
Under the assumption of Lemma 2.1, for any real value function $f\in
C^{3}(\Omega)\cap C^{2}(\overline{\Omega}),$ we have
$\displaystyle\left((\lambda_{k+2}-\lambda_{i})+(\lambda_{k+1}-\lambda_{i})\right)\int_{\Omega}|\nabla
f|^{2}u_{i}^{2}$ (2.16) $\displaystyle\leq$ $\displaystyle
2\sqrt{\left((\lambda_{k+2}-\lambda_{i})(\lambda_{k+1}-\lambda_{i})\right)\int_{\Omega}|\nabla
f|^{4}u^{2}}+\int_{\Omega}\left(2\nabla f\cdot\nabla u_{i}+u_{i}\Delta
f\right)^{2}.$
###### Proof.
Taking $g=\exp(\sqrt{-1}\alpha f),\,\alpha\in\mathbb{R}\backslash\\{0\\}$ in
(2.1), we have
$\displaystyle\alpha^{2}\left((\lambda_{k+1}-\lambda_{i})+(\lambda_{k+2}-\lambda_{i})\right)\int_{\Omega}|\nabla
f|^{2}u_{i}^{2}$ (2.17) $\displaystyle\leq$
$\displaystyle\alpha^{4}\int_{\Omega}|\nabla
f|^{4}u_{i}^{2}+\alpha^{2}\int_{\Omega}\left|2\nabla f\cdot\nabla
u_{i}+u_{i}\Delta
f\right|^{2}+(\lambda_{k+1}-\lambda_{i})(\lambda_{k+2}-\lambda_{i}).$
From (2.17), we have
$\displaystyle\left((\lambda_{k+1}-\lambda_{i})+(\lambda_{k+2}-\lambda_{i})\right)\int_{\Omega}|\nabla
f|^{2}u_{i}^{2}$ (2.18) $\displaystyle\leq$
$\displaystyle\alpha^{2}\int_{\Omega}|\nabla
f|^{4}u_{i}^{2}+\frac{1}{\alpha^{2}}(\lambda_{k+1}-\lambda_{i})(\lambda_{k+2}-\lambda_{i})+\int_{\Omega}\left|2\nabla
f\cdot\nabla u_{i}+u_{i}\Delta f\right|^{2}.$
Using the Cauchy-Schwarz inequality in (2.18), we have (2.16). ∎
###### Corollary 2.3.
Under the assumption of Lemma 2.1, for any real value function $f\in
C^{3}(\Omega)\cap C^{2}(\overline{\Omega})$ satisfying $|\nabla f|^{2}=1$, we
have
$(\lambda_{k+2}-\lambda_{k+1})^{2}\leq 16\left(\int_{\Omega}(\nabla
f\cdot\nabla u_{i})^{2}-\frac{1}{4}\int_{\Omega}(\Delta
f)^{2}u_{i}^{2}-\frac{1}{2}\int_{\Omega}(\nabla(\Delta f)\cdot\nabla
f)u_{i}^{2}\right)\lambda_{k+2}.$ (2.19)
Furthermore, we have
$\lambda_{k+2}-\lambda_{k+1}\leq
4\left(\lambda_{i}-\frac{1}{4}\int_{\Omega}(\Delta
f)^{2}u_{i}^{2}-\frac{1}{2}\int_{\Omega}(\nabla(\Delta f)\cdot\nabla
f)u_{i}^{2}\right)^{\frac{1}{2}}\sqrt{\lambda_{k+2}}.$ (2.20)
###### Proof.
From Corollary 2.2 and $|\nabla f|^{2}=1$, we have
$\left((\lambda_{k+2}-\lambda_{i})+(\lambda_{k+1}-\lambda_{i})\right)-2\sqrt{(\lambda_{k+2}-\lambda_{i})(\lambda_{k+1}-\lambda_{i})}\leq\int_{\Omega}\left(2\nabla
f\cdot\nabla u_{i}+u_{i}\Delta f\right)^{2},$
i.e.
$\left(\sqrt{\lambda_{k+2}-\lambda_{i}}-\sqrt{\lambda_{k+1}-\lambda_{i}}\right)^{2}\leq\int_{\Omega}\left(2\nabla
f\cdot\nabla u_{i}+u_{i}\Delta f\right)^{2}.$
By integration by parts, we have
$\int_{\Omega}\left(2\nabla f\cdot\nabla u_{i}+u_{i}\Delta
f\right)^{2}=4\int_{\Omega}(\nabla f\cdot\nabla
u_{i})^{2}-\int_{\Omega}(\Delta f)^{2}u_{i}^{2}-2\int_{\Omega}(\nabla(\Delta
f)\cdot\nabla f)u_{i}^{2}.$
Hence, we have
$\left(\sqrt{\lambda_{k+2}-\lambda_{i}}-\sqrt{\lambda_{k+1}-\lambda_{i}}\right)^{2}\leq
4\int_{\Omega}(\nabla f\cdot\nabla u_{i})^{2}-\int_{\Omega}(\Delta
f)^{2}u_{i}^{2}-2\int_{\Omega}(\nabla(\Delta f)\cdot\nabla f)u_{i}^{2}$ (2.21)
Multiplying (2.21) by
$\left(\sqrt{\lambda_{k+2}-\lambda_{i}}+\sqrt{\lambda_{k+1}-\lambda_{i}}\right)^{2}$
on both sides, we can get
$\displaystyle(\lambda_{k+2}-\lambda_{k+1})^{2}\leq$ $\displaystyle
4\left(\int_{\Omega}(\nabla f\cdot\nabla
u_{i})^{2}-\frac{1}{4}\int_{\Omega}(\Delta
f)^{2}u_{i}^{2}-\frac{1}{2}\int_{\Omega}(\nabla(\Delta f)\cdot\nabla
f)u_{i}^{2}\right)$
$\displaystyle\times\left(\sqrt{\lambda_{k+2}-\lambda_{i}}+\sqrt{\lambda_{k+1}-\lambda_{i}}\right)^{2}$
$\displaystyle\leq$ $\displaystyle 16\left(\int_{\Omega}(\nabla f\cdot\nabla
u_{i})^{2}-\frac{1}{4}\int_{\Omega}(\Delta
f)^{2}u_{i}^{2}-\frac{1}{2}\int_{\Omega}(\nabla(\Delta f)\cdot\nabla
f)u_{i}^{2}\right)\lambda_{k+2},$
which is the inequality (2.19).
From the Cauchy-Schwarz inequality and integration by parts, we obtain
$(\lambda_{k+2}-\lambda_{k+1})^{2}\leq
16\left(\lambda_{i}-\frac{1}{4}\int_{\Omega}(\Delta
f)^{2}u_{i}^{2}-\frac{1}{2}\int_{\Omega}(\nabla(\Delta f)\cdot\nabla
f)u_{i}^{2}\right)\lambda_{k+2}.$
Finally we have (2.20).
∎
###### Proof of Theorem 1.1 .
Let $x_{1},\;x_{2},\cdots,x_{n}$ be the standard coordinate functions in
$\mathbb{R}^{n}$. Since $|\nabla x_{l}|=1,l=1,\cdots,n$, we can use Lemma 2.3.
Taking
$f=x_{l},\,l=1,\cdots,n,\,\mbox{and}\,\,i=1$
in (2.19) and then taking sum over $l$ from $1$ to $n$, we have
$\displaystyle n(\lambda_{k+2}-\lambda_{k+1})^{2}\leq$ $\displaystyle
16\lambda_{k+2}\int_{\Omega}\sum_{l=1}^{n}\left(\frac{\partial u_{1}}{\partial
x_{l}}\right)^{2}$ (2.22) $\displaystyle=$ $\displaystyle
16\lambda_{1}\lambda_{k+2}.$
From Theorem 3.1 in [18] (see also (1.9)), from (2.22), we deduce
$\displaystyle\lambda_{k+2}-\lambda_{k+1}\leq$ $\displaystyle
4\sqrt{\frac{\lambda_{1}}{n}}\sqrt{\lambda_{k+2}}$ $\displaystyle\leq$
$\displaystyle 4\lambda_{1}\sqrt{\frac{C_{0}(n)}{n}}(k+1)^{\frac{1}{2}}$
$\displaystyle=$ $\displaystyle C_{n,\Omega}(k+1)^{\frac{1}{2}},$
where $C_{n,\Omega}=4\lambda_{1}\sqrt{\frac{C_{0}(n)}{n}}$, $C_{0}(n)$ is
given by $(\ref{diffyong})$. Since $k$ is arbitrary, this completes the proof
of Theorem 1.1. ∎
Although Corollary 1.2 can be deduced directly by taking $a=b=1$ in Corollary
1.3, its proof is interesting independently. Here we give its proof for the
upper half-plane model of hyperbolic space.
###### Proof of Corollary 1.2.
For convenience, we will use the upper half-plane model of the hyperbolic
space, that is,
$\mathbb{H}^{n}(-1)=\\{(x_{1},\cdots,x_{n})\in\mathbb{R}^{n}|x_{n}>0\\}$
with the standard metric
$\mathrm{d}s^{2}=\frac{(\mathrm{d}x_{1})^{2}+\cdots+(\mathrm{d}x_{n})^{2}}{(x_{n})^{2}}.$
Taking $r=\log x_{n},$ we have
$\mathrm{d}s^{2}=(\mathrm{d}r)^{2}+\mathrm{e}^{-2r}\sum_{i=1}^{n-1}(\mathrm{d}x_{i})^{2}.$
Since $|\nabla r|=1,\;\Delta r=-(n-1)$, taking $f=r\,\,\mbox{and}\,i=1$ in
(2.20), we have
$\displaystyle\lambda_{k+2}-\lambda_{k+1}\leq$ $\displaystyle
4\left(\lambda_{1}-\frac{1}{4}\int_{\Omega}(\Delta
r)^{2}u_{i}^{2}-\frac{1}{2}\int_{\Omega}(\nabla(\Delta r)\cdot\nabla
r)u_{1}^{2}\right)^{\frac{1}{2}}\sqrt{\lambda_{k+2}}$ (2.23) $\displaystyle=$
$\displaystyle
4\left(\lambda_{1}-\frac{(n-1)^{2}}{4}\right)^{\frac{1}{2}}\sqrt{\lambda_{k+2}}.$
By the result in [14] (see also (1.12)), from (2.23), we have
$\displaystyle\lambda_{k+2}-\lambda_{k+1}\leq$ $\displaystyle
4\left(\lambda_{1}-\frac{(n-1)^{2}}{4}\right)^{\frac{1}{2}}\sqrt{C_{0}(n)\left(\lambda_{1}+\frac{n^{2}}{4}H_{0}^{2}\right)}(k+1)^{\frac{1}{n}}$
$\displaystyle=$ $\displaystyle C_{n,\Omega}(k+1)^{\frac{1}{n}},$
where $C_{n,\Omega}$ is defined by (1.18). Since this inequality holds for any
$k$, we can deduce (1.17). ∎
## 3\. Proof of Corollary 1.3
###### Proof of Corollary 1.3.
Assume that $\Omega$ is a bounded domain in an $n$-dimensional complete
noncompact Riemannian manifold $(M,\,g)$ with sectional curvature $Sec$
satisfying $-a^{2}\leq Sec\leq-b^{2}$, where $0\leq b\leq a$ are constants.
For $p\notin\overline{\Omega}$ fixed, define the distance function by
$\rho(x)=\mbox{distance}(x,\,p)$. From Proposition 2.2 in P.15 of [31], and
$|\nabla\rho|=1,$ we have
$\nabla\rho\cdot\nabla(\Delta\rho)=-|\text{Hess}\
\rho|^{2}-\text{Ric}(\nabla\rho,\nabla\rho).$ (3.1)
Let $0\leq h\leq h_{1},\cdots,h_{n-1}\leq H$ be the eigenvalues of the
$\mathrm{Hess}\rho$. Then we have
$\displaystyle 2|\mbox{Hess}\rho|^{2}-(\Delta\rho)^{2}$ (3.2) $\displaystyle=$
$\displaystyle
2\sum_{i=1}^{n-1}h_{i}^{2}-\left(\sum_{i=1}^{n-1}h_{i}\right)^{2}$
$\displaystyle=$ $\displaystyle\sum_{i=1}^{n-1}h_{i}^{2}-\sum_{i\neq
j}h_{i}h_{j}$ $\displaystyle\leq$ $\displaystyle
h_{n-1}^{2}+h_{1}h_{2}+\cdots+h_{n-2}h_{n-1}-\sum_{i\neq j}h_{i}h_{j}$
$\displaystyle=$ $\displaystyle h_{n-1}^{2}-h_{1}h_{2}-\cdots-
h_{n-2}h_{n-1}-\sum_{i\neq j\atop i,j\leq n-2}h_{i}h_{j}$ $\displaystyle\leq$
$\displaystyle H^{2}-(n-2)^{2}h^{2}.$
From the Hessian comparison theorem (cf. [36]), under the conditions in
Corollary 1.3, we have
$a\frac{\cosh a\rho}{\sinh a\rho}\geq h_{n-1}\geq\cdots\geq h_{1}\geq
b\frac{\cosh b\rho}{\sinh b\rho}.$ (3.3)
Since $n\geq 3$ and $\frac{a^{2}}{\sinh^{2}a\rho}$ is a decreasing function of
$a$, from (3.2) and (3.3), under the conditions in Corollary 1.3, we have
$\displaystyle
2|\mbox{Hess}\rho|^{2}+2\mbox{Ric}(\nabla\rho,\,\nabla\rho)-(\Delta\rho)^{2}$
(3.4) $\displaystyle\leq$ $\displaystyle
a^{2}\frac{\cosh^{2}a\rho}{\sinh^{2}a\rho}-(n-2)^{2}b^{2}\frac{\cosh^{2}b\rho}{\sinh^{2}b\rho}-2(n-1)b^{2}$
$\displaystyle=$ $\displaystyle
a^{2}+\frac{a^{2}}{\sinh^{2}a\rho}-(n-2)^{2}b^{2}-(n-2)^{2}\frac{b^{2}}{\sinh^{2}b\rho}-2(n-1)b^{2}$
$\displaystyle\leq$
$\displaystyle-(n-1)^{2}b^{2}+(a^{2}-b^{2})+\frac{b^{2}}{\sinh^{2}b\rho}-(n-2)^{2}\frac{b^{2}}{\sinh^{2}b\rho}$
$\displaystyle\leq$ $\displaystyle-(n-1)^{2}b^{2}+(a^{2}-b^{2}).$
Taking $f=\rho\,\,\mbox{and}\,i=1$ in (2.20), we have
$\lambda_{k+2}-\lambda_{k+1}\leq
4\left(\lambda_{1}-\frac{1}{4}\int_{\Omega}(\Delta\rho)^{2}u_{1}^{2}-\frac{1}{2}\int_{\Omega}(\nabla(\Delta\rho)\cdot\nabla\rho)u_{1}^{2}\right)^{\frac{1}{2}}\sqrt{\lambda_{k+2}}.$
(3.5)
From (3.1) and (3.4), we obtain
$\displaystyle\lambda_{1}-\frac{1}{4}\int_{\Omega}(\Delta\rho)^{2}u_{1}^{2}-\frac{1}{2}\int_{\Omega}(\nabla(\Delta\rho)\cdot\nabla\rho)u_{1}^{2}$
(3.6) $\displaystyle=$
$\displaystyle\lambda_{1}+\frac{1}{4}\int_{\Omega}\left(2|\mathrm{Hess}\rho|^{2}+2\mathrm{Ric}(\nabla\rho,\,\nabla\rho)-(\Delta\rho)^{2}\right)u_{1}^{2}$
$\displaystyle\leq$
$\displaystyle\lambda_{1}-\frac{(n-1)^{2}}{4}b^{2}+\frac{a^{2}-b^{2}}{4}.$
By the result in [14] (see also (1.12)), from (3.5) and (3.6), we have
$\displaystyle\lambda_{k+2}-\lambda_{k+1}\leq$ $\displaystyle
4\left(\lambda_{1}-\frac{(n-1)^{2}}{4}b^{2}+\frac{a^{2}-b^{2}}{4}\right)^{\frac{1}{2}}\sqrt{C_{0}(n)\left(\lambda_{1}+\frac{n^{2}}{4}H_{0}^{2}\right)}(k+1)^{\frac{1}{n}}$
$\displaystyle\leq$ $\displaystyle C_{n,\Omega}(k+1)^{\frac{1}{n}},$
where $C_{n,\Omega}$ is defined by (1.20). Since this inequality holds for any
$k$, we can deduce (1.19). ∎
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* [5] M. S. Ashbaugh and R. D. Benguria, More bounds on eigenvalue ratios for Dirichlet Laplacians in $n$ dimension, SIAM J. Math Anal., 24 (1993), 1622-1651.
* [6] M. S. Ashbaugh and R. D. Benguria, Isoperimetric bounds for higher eigenvalue ratios for the $n$-dimensional fixed membrane problem, Proc. Royal Soc. Edinburgh, 126A (1993), 977-985.
* [7] M. S. Ashbaugh and R. D. Benguria, _Bounds for ratios of the first, second, and third membrane eigenvalues_ , Nonlinear Problems in Applied Mathematica, in Honor of Ivar Stakgold on this Seventieth Birthday, T. S. Angell, L. Pamela Cook, R. E. Kleinman, and W. E. Olmstead, editors, Society for Industrial and Applied Mathematics, Philadelphia, Pennsylvania, 1996, 30-42.
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Daguang Chen
Department of Mathematical Sciences, Tsinghua University, Beijing, 100084, P.
R. China
E-mail: [email protected]
Tao Zheng
Department of Mathematics, Beijing Institute of Technology, Beijing 100081, P.
R. China
E-mail: [email protected]
Hongcang Yang
Hua Loo-Keng Key Laboratory of Mathematics, Chinese Academy of Sciences,
Beijing 100080, P. R. China
E-mail:[email protected]
|
arxiv-papers
| 2013-09-28T10:19:17 |
2024-09-04T02:49:51.625981
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Daguang Chen, Tao Zheng and Hongcang Yang",
"submitter": "Tao Zheng",
"url": "https://arxiv.org/abs/1309.7446"
}
|
1309.7465
|
# Fractional quantum Hall states in charge-imbalanced bilayer systems
N. Thiebaut1 N. Regnault2,3 and M.O. Goerbig1 1 Laboratoire de Physique des
Solides, CNRS UMR 8502, Univ. Paris-Sud, F-91405 Orsay cedex, France 2
Laboratoire Pierre Aigrain, ENS and CNRS, 24 rue Lhomond, 75005 Paris, France
3 Department of Physics, Princeton University, Princeton, NJ 08544
[email protected]
###### Abstract
We study the fractional quantum Hall effect in a bilayer with charge-
distribution imbalance induced, for instance, by a bias gate voltage. The
bilayer can either be intrinsic or it can be formed spontaneously in wide
quantum wells, due to the Coulomb repulsion between electrons. We focus on
fractional quantum Hall effect in asymmetric bilayer systems at filling factor
$\nu=4/11$ and show that an asymmetric Halperin-like trial wavefunction gives
a valid description of the ground state of the system.
## 1 Introduction
In a two-dimensional electron gas subjected to a high magnetic field
repulsively interacting electrons may collectively form a gapped state
insensitive to small perturbations. Such states are called fractional quantum
Hall states (FQHS) and show up for partial (fractional) fillings of the
highest partially occupied Landau level, i.e. for rational values of the
filling factor $\nu$, which is the ratio of the number of electrons $N$ and
the number of magnetic flux quanta $N_{\phi}$ threading the system. Since the
discovery of the fractional quantum Hall effect (FQHE) in 1982 [1] new
experimental setups nourished our understanding of this multifaceted
phenomenon.
The bilayer system, which consists of two parallel 2D electron systems in a
semiconductor heterostructure exhibits FQHE at even-denominator filling
factors (such as $\nu=1/2$ [2]) is one of them. Though it is more common to
consider symmetric layers, for which charges are equally spread between the
two layers ($\nu_{\text{tot}}=2/n$), we are interested in a bilayer system
with charge imbalance, which is obtained by the application of an external
gate voltage (fig. 1). The strength of correlations between electrons that
belong to different layers depends on their distance $d$ and the energy
difference between the two layers $\Delta$ is proportional to the external
gate voltage $V_{G}$.
Beside its own richness, the FQHE in a bilayer system provides a model for its
counterpart in a wide quantum well [3, 4, 5], where even-denominator states
have also been observed [3, 4, 6]. In the latter, density functional theory
calculations indicate that the Coulomb repulsion favors an electron
concentration at the borders of the quantum well, which may thus be viewed as
a spontaneously formed bilayer system [7]. Charge-imbalanced bilayers can thus
depict wide quantum wells in which electrons split in two layers and populate
one of them preferentially, either spontaneously or under the influence of
external gate electrodes.
Here, we investigate the conditions for the existence of FQHS (at $\nu=4/11$)
in an imbalanced bilayer system by means of exact diagonalization. We
introduce asymmetric Halperin states and discuss their relevance for the
description of the ground state of the system.
Figure 1: Schematic view of the effect of a gate voltage on a bilayer system.
(Left) Setup. (Right) Effect of the gate voltage on the confinement potential.
We consider the limit where the width $w$ of each quantum well is small as
compared to the layer separation $d$, $w\ll d$.
## 2 Model
We restrict the electron dynamics to the lowest Landau level, which amounts to
considering the limit in which adjacent Landau levels are separated by a large
cyclotron gap, as compared to the characteristic Coulomb energy scale,
$e^{2}/4\pi l_{B}$. FQHS are often fully polarized due to the Coulomb
repulsion between electrons in a flat Landau energy level, and we consider
such states, that is we neglect the electron spin from now on. The effective
Hamiltonian of the bilayer system reads
$\hat{H}=\frac{1}{2}\sum_{\\{\sigma_{i}\\}}\sum_{\\{m_{i}\\}}V_{\\{m_{i}\\}}^{\\{\sigma_{i}\\}}c_{m_{1},\sigma_{1}}^{\dagger}c_{m_{2},\sigma_{2}}^{\dagger}c_{m_{4},\sigma_{4}}c_{m_{3},\sigma_{3}}+\Delta\sum_{m}(c_{m,\downarrow}^{\dagger}c_{m,{\downarrow}}-c_{m,\uparrow}^{\dagger}c_{m,{\uparrow}})\quad,$
(1)
where $\sigma=\uparrow,\downarrow$ is the layer index and $m_{i}$ the angular
momentum which labels the single-electron states in the symmetric gauge111Here
we assume that the confinement potential is sufficiently narrow (compared to
the layer separation, $w\ll d$) to avoid electrons from populating the first
excited state of the confinement potential. This permits to neglect the
wavefunction extent in the confinement direction, which assumption simplifies
the single particle Hilbert space as the layer index does not couple to the
orbital motion. [8]; $\Delta$ is the energy difference between the two layers,
controlled by the gate voltage $V_{G}$ in fig. 1. The first term on the right
hand side of eq. 1 is the usual interaction term, while the second one is a
Zeeman-like term proportional to the density imbalance. In real space the
interaction potential depends on the layer indices
$V^{\\{\sigma_{i}\\}}(r)=\left\\{\begin{array}[]{cl}\frac{e^{2}}{4\pi
r}&\mbox{ if electrons belong to the same layer (i.e. if
}\sigma_{1}=\sigma_{2}=\sigma_{3}=\sigma_{4}\text{)}\\\ &\\\
\frac{e^{2}}{4\pi\sqrt{r^{\,2}+d^{\,2}}}&\mbox{ if electrons belong to
different layers (i.e. if
}\sigma_{1}=\sigma_{3}\neq\sigma_{2}=\sigma_{4}\text{)}\\\ &\\\
0&\mbox{elsewise}\end{array}\right.$ (2)
## 3 Halperin’s wavefunctions
A large variety of FQHS in bilayer systems may be described with the help of a
generalization of Laughlin’s wavefunctions [9] to the case of electrons with
internal degrees of freedom, initially introduced by Halperin [10]. Those
wavefunctions have three integer parameters that characterise the interaction
strength between electrons in the $\uparrow$ and $\downarrow$ layers
($m_{\uparrow}$ and $m_{\downarrow}$), and electrons in different layers
($n$). In terms of the complex positions
$z_{i}^{\sigma}=x_{i}^{\sigma}-iy_{i}^{\sigma}$, for a system with
$N_{\sigma}$ electrons in the $\sigma$-layer, Halperin’s wavefunctions read
$\varPsi_{(m_{\uparrow},m_{\downarrow},n)}\big{(}\\{z_{i}^{\uparrow}\\},\\{z_{i}^{\downarrow}\\}\big{)}=\prod_{i<j\leq
N_{\uparrow}}(z_{i}^{\uparrow}-z_{j}^{\uparrow})^{m_{\uparrow}}\prod_{i<j\leq
N_{\downarrow}}(z_{i}^{\downarrow}-z_{j}^{\downarrow})^{m_{\downarrow}}\prod_{i\leq
N_{\uparrow},\,j\leq
N_{\downarrow}}(z_{i}^{\uparrow}-z_{j}^{\downarrow})^{n}\quad,$ (3)
where the normalization and usual Gaussian factors have been absorbed in a
redefinition of the Hilbert space measure
$d\mu(\\{z_{i}\\})=\prod_{i}dx_{i}dy_{i}\exp(-|z_{i}|^{2}/2l_{B}^{2})/(2\pi
l_{B}^{2})$, following Girvin and Jach [11]. Alternatively, we denote the
state represented by the Halperin wavefunction (3) the
$(m_{\uparrow},m_{\downarrow},n)$ state. Notice that not all Halperin state
(3) are capable of describing physically possible FQHS – it has been shown in
Ref. [12] within a generalization of Laughlin’s plasma analogy [9] that the
exponents must satisfy the condition
$m_{\uparrow}m_{\downarrow}\geq n^{2}$ (4)
in order to avoid phase separation of the electron species.
To compute the filling factor associated with a particular Halperin state, one
uses the fact that the maximum value of individual orbital momenta is given by
the number of flux quanta $N_{\phi}$ threading the system. The polynomial
expansion of eq. (3) yields a maximal exponent
$(N_{\sigma}-1)\;m_{\sigma}+N_{-\sigma}\;n$ for the position of the $i$-th
electron in the $\sigma$-layer that must equal $N_{\phi}$. This can be written
in matrix form as
$\begin{pmatrix}N_{\phi}+m_{\uparrow}\\\
N_{\phi}+m_{\downarrow}\end{pmatrix}=\begin{pmatrix}m_{\uparrow}&n\\\
n&m_{-}{\downarrow}\end{pmatrix}\begin{pmatrix}N_{\uparrow}\\\
N_{\downarrow}\end{pmatrix}\quad,$
and after inverting the matrix, one deduces the relationship between the total
number of electrons $N$ and $N_{\phi}$,
$N=N_{\uparrow}+N_{\downarrow}=\frac{1}{m_{\uparrow}m_{\downarrow}-n^{2}}\Big{[}(m_{\uparrow}+m_{\downarrow}-2n)N_{\phi}+2m_{\uparrow}m_{\downarrow}-n(m_{\uparrow}+m_{\downarrow})\Big{]}\quad.$
(5)
The filling factor $\nu$ is defined in the thermodynamic limit, for which we
obtain
$\nu=\lim_{N,N_{\phi}\rightarrow\infty}\frac{N}{N_{\phi}}=\frac{m_{\uparrow}+m_{\downarrow}-2n}{m_{\uparrow}m_{\downarrow}-n^{2}}\quad.$
(6)
We define the polarization as the population difference between the layers
$P_{z}=\frac{S_{z}}{N/2}$ (7)
in terms of the $z$-component of the total layer pseudospin
$S_{z}=\frac{N_{\uparrow}-N_{\downarrow}}{2}=\frac{(m_{\downarrow}-m_{\uparrow})(N_{\phi}-n)}{2(m_{\uparrow}m_{\downarrow}-n^{2})}\quad,$
(8)
which is zero for symmetric states (i.e. if $m_{\uparrow}=m_{\downarrow}$).
Now we focus on asymmetric Halperin states. More precisely we are interested
in the $(5,3,2)$ state, which satisfies the condition (4) since it is the
simplest (i.e. lowest correlation factors) asymmetric Halperin state for which
the filling factor $\nu_{532}=4/11$ does not belong to the usual composite-
fermions sequence $\nu_{\text{CF}}=p/(2sp\pm 1)$ [13]. It may thus describe
states the existence of which relies on layer asymmetry. Nevertheless it
should be mentioned that $\nu=4/11$ is also the filling factor of a second-
generation composite fermion state in monolayer systems [14, 15, 16, 17]
generated by interacting composite fermions that has been proposed as a
possible explanation for an experimentally observed FQHS at this filling [14].
## 4 Exact diagonalization
In order to check the physical relevance of the $(5,3,2)$ state we calculate,
with the help of exact diagonalization [18], two quantities for the Coulomb
interaction (2) in the lowest Landau level: the total angular momentum of the
ground state and its polarization as a function of the distance $d$ and of the
gate-induced energy difference $\Delta$ between the two layers. The
calculations are perfomed on the sphere geometry [19]. For the $(5,3,2)$
state, the relation between $N_{\Phi}$ and $N$ is given by
$N_{\Phi}=\frac{11}{4}N-\frac{7}{2}$. Finite-size systems induce a constraint
on the number of electrons [see eq. (5)]. Only sizes of the form $N=2+4n$,
$n\in\mathbb{N}$ can satisfy this constraint.
Here we will focus on $N=6$ and $N=10$ which are the only accessible sizes.
FQHS correspond to values of $d$ and $\Delta$ for which the gap is finite and
the total angular momentum equals zero. Finally we compute the overlap between
the $(5,3,2)$ state and the exact ground state for values of $d$ and $\Delta$
which correspond to a FQHS.
The polarization of the ground state is shown in fig. 2. For sufficiently high
values of $\Delta$ the system is fully polarized, i.e all electrons reside in
a single layer and one thus obtains the monolayer system with $\nu=4/11$,
studied in Refs. [15, 16, 17]. Conversely for large distances $d$ the
polarization tends to be zero. In this situation the potential
$V(r)=(r^{\,2}+d^{\,2})^{-1/2}$ between two electrons in different layers is
small and the system behaves as two uncorrelated layers with individual
filling factors $\nu=2/11$. It is likely that for such a low filling factor
the system tends to form a Wigner crystal [20, 21, 22], and FQHS may therefore
be ruled out.
(a) N=6: $P_{z}^{(532)}=-1/3$ (b) N=10: $P_{z}^{(532)}=-2/5$
Figure 2: Polarization of the bilayer system as a function of the distance
$d/l_{B}$ and energy shift $\Delta$ (energies are given in units of
$e^{2}/4\pi l_{B}$).
The polarization of the $(5,3,2)$ state can be computed directly from eq. (7)
and eq. (8)
$P_{z}^{(532)}=-\frac{1}{2}+\frac{1}{N}$ (9)
and may be compared with that obtained from exact diagonalization for $N=6$
and 10 particles, as a function of the layer separation $d/l_{B}$ and
$\Delta/(e^{2}/4\pi l_{B})$ (see fig. 2). Indeed, one finds that in
intermediate regimes the polarization of the ground state is the one of the
$(5,3,2)$ state (orange regions).
We now turn to the characterization of the exact ground state. Figure 3 shows
the total angular momentum $L$ of the ground state. This latest is
rotationally invariant when $L=0$, as required for an incompressible state. It
is zero in the region where the ground-state polarization matches that of the
$(5,3,2)$ state, except for large layer separations. As already mentioned in
the discussion above, for large distances inter-layer correlations are too
weak to allow the system to be in an incompressible state, and an
inhomogeneous density state is favored, such as one would for example expect
for a Wigner crystal ($d\gtrsim 3.8l_{B}$).
Figure 3: Total orbital momentum of the bilayer system for $N=6$ electrons as
a function of the distance $d/l_{B}$ and energy shift $\Delta$, in units of
$e^{2}/4\pi l_{B}$.
The results presented in figs. 2 and 3 indicate that in a certain range of
parameters $d/l_{B}$ and $\Delta/(e^{2}/4\pi l_{B})$, exact diagonalization
provides a ground state that matches some physical properties of the $(5,3,2)$
state. To corroborate the relevance of this state, we have calculated the
overlap between the exact ground state obtained for the interaction (2) and
the $(5,3,2)$ state. In order to obtain the $(5,3,2)$ state numerically, we
have used exact diagonalization for a model interaction given in terms of the
appropriate pseudopotentials [19]. The results are presented in fig. 4 as a
function of $d/l_{B}$. The overlap turns out to be high ($>85\%$) for
$d\lesssim 2l_{B}$, such that in regions of matching polarizations and for
sufficiently short interlayer distances the $(5,3,2)$ state gives a good
description of the ground state.
Figure 4: Overlap of the $(5,3,2)$ state with the exact ground state as a
function of $d/l_{B}$, in the corresponding polarization sector (9).
## 5 Conclusion
We have performed an exact-diagonalization study of the bilayer system for
various distances between the layers and energy shifts, i.e. the energy
difference between electrons in the two layers. The ground state is
incompressible over a certain range of parameters, which corresponds to an
interlayer distance of roughly three times the magnetic length. Asymmetry
between layers does not necessarily destroy the FQHS, and the FQHE at exotic
fractions may thus be observed in an asymmetric bilayer system. Asymmetric
Halperin’s wavefunctions are good candidates for the accurate description of
those states. Also, using the bilayer modeling of a wide quantum well, a
similar behavior could be expected in asymmetric wide quantum wells, the
asymmetry being induced, for instance, by a back gate voltage that enables one
to tune the electronic density of the system. Experimentalists reported a FQHE
at $\nu=4/11$ in a wide quantum well ($w\sim 3\;l_{B}$) in 2003 [14]. Since
this filling factor corresponds to the one of the $(5,3,2)$ state, one may
wonder about its relevance for the explanation of this experimental
observation. Within this picture this FQHS would be induced by the shape of
the confinement potential, in contrast to the picture of second-generation
composite fermions that have been proposed as an alternative explanation for
the 4/11 state [14, 15, 16, 17].
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|
arxiv-papers
| 2013-09-28T14:45:52 |
2024-09-04T02:49:51.634190
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "N. Thiebaut, N. Regnault and M.O. Goerbig",
"submitter": "Nicolas Thi\\'ebaut",
"url": "https://arxiv.org/abs/1309.7465"
}
|
1309.7473
|
# A fast map-making preconditioner for regular scanning patterns
Sigurd K. Næss [email protected] Thibaut Louis
[email protected] Sub-department of Astrophysics, University of
Oxford, Keble Road, Oxford, OX1 3RH, UK
###### Abstract
High-resolution Maximum Likelihood map-making of the Cosmic Microwave
Background is usually performed using Conjugate Gradients with a
preconditioner that ignores noise correlations. We here present a new
preconditioner that approximates the map noise covariance as circulant, and
show that this results in a speedup of up to 400% for a realistic scanning
pattern from the Atacama Cosmology Telescope. The improvement is especially
large for polarized maps.
## I Introduction
As the resolution and sensitivity of Cosmic Microwave Background (CMB)
experiments increase, so do the computational resources needed to analyze
their data. Because modern detectors are background-limited, the only way to
significantly increase sensitivity is to increase the number of detectors. The
last decades have seen an increase from tens of detectors to thousands of
detectors in experiments like ACT Dünner et al. (2013); Niemack et al. (2010),
SPT Carlstrom et al. (2011); Austermann et al. (2012), POLARBEAR Kermish et
al. (2012) and Keck Kernasovskiy et al. (2012), and plans already exist for
experiments with $10^{5}-10^{6}$ detectors Abazajian et al. (2013). Reducing
the data from all these detectors into a coherent map of the sky presents a
significant computational challenge, and already with 1000-detector-class
experiments this step is the most important bottleneck of the data analysis
pipeline Dünner et al. (2013). It is therefore important to investigate ways
to speed up this process.
Three main classes of map-makers are in popular use: _Maximum likelihood map-
makers_ Tegmark (1997); Cantalupo et al. (2010); de Gasperis et al. (2005);
Traficante et al. (2011); Dünner et al. (2013); QUIET Collaboration et al.
(2012), which are slow, but produce unbiased, optimally noise-weighted maps;
faster but slightly less accurate _destripers_ Keihänen et al. (2005); Sutton
et al. (2010); and biased and sub-optimal but very fast _naive map-makers_
QUIET Collaboration et al. (2012); Schaffer et al. (2011). The topic of this
paper is a method for significantly speeding up maximum likelihood map-makers.
Assuming a linear detector response, we can model the time-ordered data $d$
via the the linear system
$\displaystyle d=Pm+n,$ (1)
where $m$ is the pixelized map of the sky, the pointing matrix $P$ is a
sparse111The pointing matrix will be sparse if we solve for a beam-convolved
map. For variable beams or asymmetric beams, one may want to reconvolve to a
standard beam as part of map-making. This can be done using $P$, at the cost
of some of its sparsity. We do not consider this case here. matrix mapping
from pixels to samples, and $n$ is the time-domain noise, which we assume to
be gaussian with covariance matrix $N$. The maximum likelihood solution for
$m$ is given by the map-making equation Tegmark (1997),
$\displaystyle(P^{T}N^{-1}P)m$ $\displaystyle=P^{T}N^{-1}d.$ (2)
This is of the form $Ax=b$, and while the matrices involved are usually too
large to solve by direct inversion, the system is amenable to solution by
Preconditioned Conjugate Gradients (PCG) Press et al. (2007) provided a good
preconditioner can be found222Without a preconditioner, the number of
iterations needed for conjugate gradients is proportional to the condition
number of the matrix $A$. By applying a preconditioner $M$, one is effectively
solving the system $MAx=Mb$. The goal is then to choose $M$ such that $MA$ is
as well-conditioned as possible. This can be acchevied if $M\approx A^{-1}$..
The most commonly used map-making preconditioners are the binned Natoli et al.
(2001); Ashdown et al. (2007); Cantalupo et al. (2010) and Jacobi Doré et al.
(2001); de Gasperis et al. (2005); Cantalupo et al. (2010); Dünner et al.
(2013) preconditioners. The binned preconditioner $M_{\textrm{B}}$
approximates the time covariance matrix $N$ as diagonal (i.e. it ignores
correlations), which for pointing matrices where only one pixel is hit per
sample results in a diagonal pixel-space covariance matrix.
$\displaystyle N^{-1}_{ij}$ $\displaystyle\approx
N^{-1}_{ii}\delta_{ij}\textrm{ and }P_{ti}P_{tj}\propto\delta_{ij}\Rightarrow$
(3) $\displaystyle A_{ij}$ $\displaystyle\equiv
P_{ti}N^{-1}_{tt^{\prime}}P_{t^{\prime}j}\approx
P_{ti}^{2}N^{-1}_{tt}\delta_{ij}\equiv{M_{\textrm{B}}}^{-1}_{ij}$ (4)
The Jacobi preconditioner $M_{\textrm{J}}$ simplifies one step further, and
assumes that every detector has the same variance, $N=aI$, resulting in
$\displaystyle A_{ij}$ $\displaystyle\approx
aP_{ti}^{2}\delta_{ij}\equiv{M_{\textrm{J}}}^{-1}_{ij}.$ (5)
The proportionality factor $a$ is usually set to 1, as PCG is insensitive to
an overall scaling of the preconditioner.
$-10^{-4}$ $10^{-4}$
Figure 1: A $4.2^{\circ}$ by $3.2^{\circ}$ subset of a row from the pixel-
space inverse correlation matrix from a patch from ACT (each row in the matrix
corresponds to a two-dimensional map). In order to highlight the correlation
structure, the color scale is capped at $\pm 10^{-4}$. The significantly
correlated area has a quite complicated shape, which is driven by the scanning
pattern and focal plane layout.
$-10^{-3}$ $10^{-3}$
Figure 2: A mosaic of the local correlation structure for a $12^{\circ}$ by
$9^{\circ}$ subset of an ACT patch. A set of evenly spaced pixels were chosen
as reference points, and for each the correlation in a neighborhood around it
was calculated. Each pixel in the map is colored according to its correlation
relative to the nearest reference point. The correlation structure is very
uniform in the whole central region of the map.
The assumption of independent noise is quite inaccurate. All realistic
experiments have at least some time-correlation in the noise, and usually also
correlations between different detectors. Additionally, filters will also
generally introduce correlations.
While computing the full, exact $A$ is often too expensive, single rows of it
can be computed at the same expense as one CG iteration:
$\displaystyle A_{ij}$ $\displaystyle=A_{ik}\delta_{kj}=(A\vec{e}_{i})_{j}$
(6)
Here $(\vec{e}_{i})_{j}=\delta_{ij}$ is the pixel-space basis vector
corresponding to pixel $i$. An example of what such a row looks like for the
Atacama Cosmology Telescope (ACT) Dünner et al. (2013) can be seen in figure
1.
It is clear that the approximation of no correlation is quite inaccurate.
However, since the correlation structure is driven by the scanning pattern and
relative position of the detectors in the focalplane, the correlation
structure should be the same for all pixels which are scanned the same way.
ACT, which used long-duration, small-amplitude drift scans, pixels at the same
declination but different right ascension will be hit by the same phase of the
same scanning motion, and should therefore have the same correlation
structure.333 Circulant correlation is not a good approximation for every
experiment. It is suitable for constant elevation dift scans, as employed by
ACT, SPT and POLARBear, but we expect it to work poorly for full-sky scanning
patterns.
Indeed, that is what measurements show (see figure 2). It may therefore be a
good approximation to assume that every point in the map has the same relative
correlation structure, i.e. that the correlation between two points on the sky
only depends on their relative position. If this is the case, then it is
possible to choose a pixelization where the correlation only depends on the
difference between pixel numbers, and hence that the pixel correlation matrix
is _circulant_ 444 For example, if
$\textrm{corr}(\vec{x}_{1},\vec{x}_{2})=f(\vec{x}_{1}-\vec{x}_{2})$, where
$\vec{x}$ are coordinates, then a pixelization scheme $\vec{x}=G(p)$, where
$p$ is a pixel index and G is a linear function will fulfill
$\textrm{corr}(p_{1},p_{2})=f(G(p_{1}-p_{2}))$, resulting in a circulant
correlation matrix..
A circulant matrix has the nice property of being diagonal in the frequency
domain, which means that it can be computed, stored and applied cheaply, at
the cost of a few FFTs. Hence, the constant correlation approximation is
promising as a preconditioner for solving the map-making equation.
## II Implementation
The inverse pixel covariance matrix $A$ can be decomposed into variance and
correlation such that $A=\Sigma^{T}U\Sigma$. Here $\Sigma$ is diagonal (block-
diagonal in the case of polarization) in pixel space, and corresponds to a map
of the inverse standard deviation per pixel. As per the binned preconditioner,
this can be approximated as
$\displaystyle(\Sigma^{T}\Sigma)_{ij}$
$\displaystyle=P_{ti}^{2}N^{-1}_{tt}\delta_{ij}.$ (7)
The correlation matrix $U$ is in general a dense matrix, but as noted above,
it can often be approximated as circulant555The matrix will be circulant
provided that the correlation structure is position-independent, that a
constant offset in each coordinate corresponds to a constant pixel offset, and
provided that indices wrap around at the edges.. Therefore, the constant
correlation preconditioner replaces $U$ with a circulant matrix $Q$, such that
$\displaystyle A$
$\displaystyle\approx\Sigma^{T}Q\Sigma\equiv{M_{\textrm{C}}}^{-1}$ (8)
This relation can be inverted to give us an expression for $Q$ in terms of
$A$,
$\displaystyle Q=$ $\displaystyle{\Sigma^{T}}^{-1}A\Sigma^{-1}\Rightarrow$ (9)
$\displaystyle Q_{ij}=$
$\displaystyle{\Sigma^{T}}^{-1}_{ii}(A\vec{e}_{i})_{j}\Sigma^{-1}_{jj}.$ (10)
Since $Q$ is circulant, i.e. $Q_{ij}=Q_{0,j-i}=q_{j-i}$, we have
$\displaystyle(FQF^{-1})_{ff^{\prime}}$
$\displaystyle=(Fq)_{-f}\delta_{ff^{\prime}}=(Fq)^{*}_{f}\delta_{ff^{\prime}}$
(11)
using forward and backwards Fourier transforms
$F_{fj}\equiv\mathrm{e}^{-\frac{2\pi ijf}{N}}$ and
$F^{-1}_{jf}\equiv\frac{1}{N}\mathrm{e}^{\frac{2\pi ijf}{N}}$, where $N$ is
the number of rows in the matrix.
With this in hand, the preconditioner can be applied as
$\displaystyle M_{\textrm{C}}b$
$\displaystyle=\Sigma^{-1}Q^{-1}{\Sigma^{-1}}^{T}b$
$\displaystyle=\Sigma^{-1}F^{-1}{(Fq)^{*}}^{-1}F{\Sigma^{-1}}^{T}b.$ (12)
$\Sigma^{-1}$ and ${(Fq)^{*}}^{-1}$ can be precomputed, so the cost of
applying the preconditioner is simply that of two FFTs and three diagonal
matrix multiplications.
The choice of the reference pixel at which the correlation is measured is
somewhat arbitrary. We used the pixel nearest the center of the map, but other
locations not too close to the edge of the map should also work.
Figure 3: Example TOD noise power spectrum from the simulation, based on noise
behavior from ACT. Low frequencies are dominated by atmospheric noise, while
the increase at high frequencies is caused by a Butterworth filter.
## III Regularization
While constant correlation is a good approximation for relatively short-scale
correlations, it works less well for long-distance correlations, and regions
near the edges666Near the edges the telescope must decelerate in order to
reverse the scanning direction, which makes the correlation structure
different there than in the center. Applying the preconditioner as described
above to realistic cases results in the appearance of large scale modes which
change extremely slowly during the subsequent CG iteration.
A way around this is to artificially limit the range of the correlations that
are modeled, by multiplying $q$ by a Gaussian. For ACT, a standard deviation
of 20 pixels was found to be effective, but this will depend on the scanning
pattern, and some experimentation may be needed to find the optimal number.
Figure 4: The residual A-norm Strakos & Tichy (2008) as a function of
iteration number for a simple binned preconditioner and the constant
correlation preconditioner. The latter converges roughly twice as fast as the
former according to this criterion. The horizontal blue line indicates the
level at which the maps have mostly stopped changing visually.
## IV Polarization
The previous discussion assumed that each pixel only had a single degree of
freedom, e.g. temperature-only maps of the sky. In the case of polarization,
each pixel has several correlated components, typically the Stokes parameters
T, Q and U Zaldarriaga & Seljak (1997), but this can instead be expressed as a
larger number of block-correlated single-component pixels. This results in
$\Sigma$ being block-diagonal with e.g. one (T,Q,U)-block per physical pixel,
while $Fq$ becomes a vector of similar blocks. And instead of a single row of
$A$ needing to be measured, all the rows corresponding to a given physical
pixel now need to be computed (i.e. $\vec{e}_{i\alpha}$ for all components,
where Greek indices indicate polarization components). Hence, eq. (10) becomes
$\displaystyle Q_{i\alpha j\beta}=$ $\displaystyle{\Sigma^{T}}^{-1}_{i\alpha
i\gamma}(A\vec{e}_{i\gamma})_{j\delta}\Sigma^{-1}_{j\delta j\beta}$ (13)
Aside from that, everything works the same.
## V Test setup
| Step 5 | Step 15 | Step 45 | Step 115 | Step 340
---|---|---|---|---|---
Binned | | | | |
Const. corr. | | | | |
Figure 5: Example temperature maps from the CG solution process for the binned (top) and constant correlation (bottom) preconditioners. The rows correspond to steps 5, 15, 45, 115 and 340 from left to right. The steps are chosen such that the binned map in column $n$ is as similar as possible as the constant correlation map in column $n-1$. We see that the constant correlation preconditioner visually converges about 3 times faster than the binned one. The maps have been cropped for compactness of presentation. Step 15 | Step 70 | Step 250
---|---|---
| |
Figure 6: Comparison of the input (blue), binned (red) and constant
correlation (green) power spectra at CG steps 15 (left), 70 (middle) and 250
(right). Unlike the map-space plots in figure 5, where only the large-scale
convergence is visible, we can here clearly see the convergence at all scales.
Small-scale convergence is much slower for EE and BB than for TT, and form the
bottleneck for the CG solver if we ignore the $\ell<500$ modes. The constant
correlation preconditioner is 3-5 times faster than the binned one here. This
is quantified more precisely in figure 9.
We tested the preconditioners on a simulated time-ordered data (TOD) based on
the scanning pattern for 64 detectors from each of 417 15-minute scans of a
subset of ACT’s southern patch centered at
$\alpha=56^{\circ},\delta=-53^{\circ}$. Each scan was a constant elevation
drift scan with amplitude of $3.5^{\circ}$ in azimuth, and the scan centers
were spread over 10 steps in elevation, covering a patch of about $11^{\circ}$
by $8^{\circ}$ degrees. Odd steps in elevation scanned while rising and the
even ones when setting. This resulted in most pixels being hit from two
directions, and hence the x-shaped correlation pattern seen in figures 1-2. An
example of a noise power spectrum used in the simulation can be seen in figure
3.
The simulated detectors were polarization-sensitive, with each detector
measuring a linear combination $T+\cos(2\psi)Q+\sin(2\psi)U$ of the local
radiation field, with each detector having a different, randomly chosen
detector angle $\psi$. While an ACT-like noise model, including the effects of
atmosphere and inter-detector noise correlations was assumed in the map-making
step, no noise was added to the simulated TOD in order to allow the
convergence to be studied all the way to the highest multipoles777 This is
valid since the convergence rate of PCG is mostly independent of the noise
level of the right-hand side after the first few iterationsNatoli et al.
(2001). However, with higher noise, higher CG errors also become acceptable,
so the number of iterations needed for CG errors to be subdominant will be
smaller for realistic noise levels.. For the same reason, the simulated input
CMB did not include a beam, and was pixelated at the same resolution as the
output map, in order to avoid subpixel noise.
We then solved the map-making equation for this data set using PCG, first
using the binned preconditioner described in equation (4), and then the
constant correlation approximation described in this paper. Each was run for
600 CG iterations, with intermediate maps being output for every 5 steps.
## VI Results
The constant correlation preconditioner visually converges roughly 3 times
faster than the binned one, as shown for the temperature map in figure 5.
Likewise, the residual A-norm Strakos & Tichy (2008) from the CG solver (shown
in figure 4) also shows a significant improvement in convergence: roughly a
factor 2 according to this metric888The A-norm $||x||_{A}$ of a vector $x$ is
defined as $\sqrt{x^{T}Ax}$, where $A$ is $P^{T}N^{-1}P$ in our case. The
error A-norm after $i$ CG steps is $||x_{i}-x||_{A}$, where $x_{i}$ is our
estimate after $i$ steps, and $x$ is the true map. When the true $x$ is
unknown, $||x_{i}-x||_{A}$ can be estimated as
$||x_{i}-x||_{A}\approx\sum_{j=i}^{i+d-1}\gamma_{j}||r_{j}||^{2}$, where
$r_{j}$ and $\gamma_{j}$ are two internal variables in the CG algorithm at
step $j$, and $d$ is is an integer that controls the accuracy of the estimate
(4 in our case).
Figure 7: An example of the convergence of a single bin in the power spectrum,
in this case that centered on $\ell=3200$. This shows the absolute difference
between the recovered spectrum $C_{\ell}$ and the spectrum $\hat{C}_{\ell}$ of
the input map in units of cosmic variance $\sigma_{\ell}^{\textrm{cv}}$, as a
function of the CG step in cosmic variance units. Shown are curves for TT
(red), EE (green) and BB (blue) for the binned (solid) and constant
correlation (dashed) preconditioners. The trend lines are based on binned
averages of the errors for many individual CG steps. This binning reduces the
jitter, and makes the trend lines suitable for measuring the time needed to
converge to a given level.
However, neither of these tests take into account the fact that not all scales
in the map are equally interesting. To remedy this, figure 6 compares the
binned and constant correlation power spectra with that of the simulated input
map. These spectra were computed using the method described in Louis et al.
(2013). On large scales ($\ell<500$), this tells the same story as the maps
did: The larger the scale, the more slowly it converges, with the constant
correlation preconditioner being about 3 times faster than the binned one.
Somewhat surprisingly, a similar phenomenon occurs at the small scales. For
$\ell>2000$, higher $\ell$ results in slower convergence, and this is
especially prominent for the EE and BB power spectra. On all scales, however,
the constant correlation preconditioner appears to converge several times
faster than the binned one.
In order to quantify the convergence more precisely, we consider the time at
which the absolute error in a given multipole-bin reaches 0.1 times cosmic
variance in that bin. While somewhat arbitrary, this choice ensures that CG
errors are guaranteed to be sub-dominant in the power spectrum, regardless of
the noise properties of the actual experiment. Figure 7 shows the convergence
of TT, EE and BB for both preconditioners for a typical multipole-bin. The
overall trend for each component is an initial rapid fall followed by a slower
decay, with both being significantly faster for the new preconditioner,
particularly for the polarization spectra.
Figure 8: The number of CG steps needed for each multipole-bin to converge to
0.1 times cosmic variance, for each of TT (red), EE (green) and BB (blue) for
the binned (solid) and constant correlation (dashed) preconditioners.
We found that binning the errors in bins of $\sim 50$ conjugate gradients
steps and using linear interpolation between these bins resulted in a robust
estimate of when each spectrum bin reaches the convergence criterion. The
resulting convergence times can be seen in figure 8, and confirm our earlier
finding that the largest and smallest scales converge more slowly. The figure
also highlights how much trouble the binned preconditioner has with the EE and
especially BB spectra, where it performs much more poorly relative the
constant correlation preconditioner than we see for the TT spectrum. We
speculate that this is due to the X-shaped correlation structure introduced by
our scanning pattern. In binned maps, which ignore the correlations, this
introduces spurious X-shaped patterns in both Q and U, corresponding to
spurious signal in both E and B. With the constant correlation preconditioner,
these are partially corrected because some of the correlation structure is
taken into account.
Figure 9: The ratio of the convergence times to 0.1 times cosmic variance for
the binned and constant correlation preconditioners, per multipole bin for
each of TT, EE and BB. For TT, the speedup is typically between 20% and 200%,
while for BB the speedup ranges from 200% at large scales to about 400% at
small scales, with EE being intermediate. Large points here indicate cases
where the binned preconditioner did not converge to 0.1 times cosmic variance
in time. For these points, the comparison was performed at the lowest
threshold where both converged.
We summarize the performance characteristics of the new preconditioner in
figure 9, which shows its relative speed gain vs. the baseline binned
preconditioner. Depending on $\ell$, we have a speedup ranging from 20% to
200% for TT, from 150% to 300% for EE, and from 200% to 400% for BB, with the
greatest relative improvement happening at the scales that converge most
slowly.
## VII Summary
The structure of a CMB map’s pixel covariance matrix is determined by the
noise properties of the time-ordered data and the scanning pattern of the
telescope. For constant elevation drift scans like those employed by ACT, SPT
and POLARBEAR, this results in an approximately circulant covariance. We have
developed a new preconditioner for conjugate gradient solutions of the map-
making equation which exploit this property by deconvolving the correlations
in harmonic space, an operation which is very cheap due to the Fourier
representation of a circulant matrices being diagonal.
For a realistic scanning pattern and noise model the preconditioner results a
speedup of 20% to 200% for temperature and 150% to 400% for polarization
compared to a binned preconditioner.
Convergence speed might potentially be further improved by allowing the
correlation pattern to change slowly across the map, for example by
decomposing the map into overlapping tiles, and applying the constant
correlation preconditioner separately to each tile, followed by a merging
operation. Our preliminary attempts at such a tiled preconditioner have
however not been able to beat the simple constant correlation approximation
presented here.
###### Acknowledgements.
The authors would like to thank Jo Dunkley and Johannes Noller for useful
discussion and suggestions, and Jon Sievers for testing the preconditioner in
another map-maker. We also thank the ACT collaboration for access to internal
ACT data used in the simulations. Computations were performed on the gpc
supercomputer at the SciNet HPC Consortium. SciNet is funded by: the Canada
Foundation for Innovation under the auspices of Compute Canada; the Government
of Ontario; Ontario Research Fund - Research Excellence; and the University of
Toronto. SN and TL are supported by ERC grant 259505.
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|
arxiv-papers
| 2013-09-28T16:14:31 |
2024-09-04T02:49:51.640215
|
{
"license": "Public Domain",
"authors": "Sigurd N{\\ae}ss, Thibaut Louis",
"submitter": "Sigurd Kirkevold N{\\ae}ss",
"url": "https://arxiv.org/abs/1309.7473"
}
|
1309.7527
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# Structured Spectrum Allocation and User Association in Heterogeneous
Cellular Networks
Wei Bao and Ben Liang
Department of Electrical and Computer Engineering, University of Toronto,
Canada
Email: {wbao, liang}@comm.utoronto.ca
###### Abstract
We study joint spectrum allocation and user association in heterogeneous
cellular networks with multiple tiers of base stations. A stochastic geometric
approach is applied as the basis to derive the average downlink user data rate
in a closed-form expression. Then, the expression is employed as the objective
function in jointly optimizing spectrum allocation and user association, which
is of non-convex programming in nature. A computationally efficient Structured
Spectrum Allocation and User Association (SSAUA) approach is proposed, solving
the optimization problem optimally when the density of users is low, and near-
optimally with a guaranteed performance bound when the density of users is
high. A Surcharge Pricing Scheme (SPS) is also presented, such that the
designed association bias values can be achieved in Nash equilibrium.
Simulations and numerical studies are conducted to validate the accuracy and
efficiency of the proposed SSAUA approach and SPS.
## I Introduction
Traditional single-tiered macro-cellular networks provide wide coverage for
mobile user equipments (UEs), but they are insufficient to satisfy the
exploding demand driven by modern mobile traffic, such as multimedia
transmissions and cloud computing tasks. One efficient means to alleviate this
problem is to install a diverse set of small-cells (e.g., picocells and
femtocells), overlaying the macrocells, to form a multi-tiered heterogeneous
cellular network [1]. Each small-cell is equipped with a shorter-range and
lower-cost base station (BS), to provide nearby UEs with higher-quality
communication links with lower power usage.
However, in the presence of multiple tiers of BSs in a cellular network, user
association control becomes more challenging. A most direct approach is
association by maximum received power, in which UEs are associated with the BS
(in any tier) with the largest received power. However, in this case, because
small-cell BSs are transmitting at lower power levels, only the UEs very close
to them will connect with them, while most other UEs are still crowding in
macrocells, leading to degraded performance. An example is shown in Fig. 1(a),
in which many UEs are occupying the macrocells, while some small-cells are
nearly empty.
In order to resolve this issue, a flexible user association approach may be
employed [2, 3], in which each tier of BSs is assigned a user association bias
value, and a UE is associated with a BS with the largest received power
multiplied by the bias value. If small-cell BSs are assigned with larger
association bias values, the small-cells are “expanded” accordingly. This can
result in a more balanced mobile traffic pattern and thus better network
performance. Fig. 1(b) shows an example of flexible user association. However,
if the association bias values for small-cell BSs are too large, it will cause
improper expansions of small-cells such that UEs at their cell-edge may suffer
from inadequate received power. As a consequence, the association bias values
should be properly designed so that the overall network performance is
optimized.
Further complicating the resource management problem in a multi-tier cellular
network, the radio spectrum licensed by the network operator needs to be
shared by BSs of widely different power and coverage areas. How to optimally
allocate spectrum among different tiers is an important open problem. In order
to avoid cross-tier interference, and the prohibitive complexity in tracking
and provisioning for such interference especially with unplanned deployment of
small cells, a disjoint spectrum mode is commonly advocated [4, 5, 6], where
different tiers of BSs are allocated non-overlapping portions of the spectrum.
Even so, it is still a challenging problem to properly divide the spectrum for
optimal network performance.
(a) Association by maximum received power.
(b) Flexible user association.
Figure 1: An example of a three-tier cellular network. Macrocell BSs, picocell
BSs, and femtocell BSs are represented by squares, circles, and triangles
respectively; UEs are represented by dots; blue lines show cell boundaries.
In this work, our objective is to study jointly optimal spectrum allocation
and user association in a heterogeneous celluar network with multiple tiers of
BSs. First, we develop a stochastic geometric model to study the network
performance analytically. A closed-form expression for the average downlink UE
data rate is derived, which is then employed as the objective function for
jointly optimizing the spectrum allocation among tiers and the user
association bias values.
This resultant optimization problem is of non-convex programming in nature and
cannot be solved with a standard method. Instead, we explore two important
structures in solving the problem. Referred to as the density thresholding
structure, we show that the problem can be studied separately over sparse UE
and dense UE scenarios, divided by a parameter specific threshold. Referred to
as the priority ordering structure, we show that a tier with higher BS density
should have higher priority in spectrum allocation. Based on these
observations, we propose a computationally efficient Structured Spectrum
Allocation and User Association (SSAUA) approach to solve the problem
optimally in the sparse UE scenario, and near-optimally in the dense UE
scenario with a quantified performance bound.
Finally, toward practical implementation of SSAUA, we propose a Surcharge
Pricing Scheme (SPS), such that the designed association bias values can be
achieved in Nash equilibrium. Hence, each UE is incentivized to adopt the
proposed design with individual rationality.
The rest of the paper is organized as follows. In section II, we discuss the
relation between our work and prior works. In Section III, we present the
system model. In Sections IV, V, and VI, we present our contributions in UE
data rate derivation, SSAUA design, and SPS, respectively. In Section VII, we
present numerical results. Finally, conclusions are given in Section VIII.
## II Related Works
### II-A Stochastic Geometry as Analytical Basis
Stochastic geometry [7, 8, 9, 10] is a powerful mathematical modeling tool to
analyze the performance (e.g., outage probability and data rate) of cellular
networks with random spatial patterns of UEs and BSs. In this work, we focus
on the downlink user data rate as performance measure. Pioneering works on
downlink performance analysis using stochastic geometry include [11], [12],
and [13], for either the single-tier or the multi-tier case. None of them
considered spectrum allocation or user association.
### II-B Spectrum Allocation and User Association
Most prior studies considered either spectrum allocation or user association
separately. For example, [6, 14] studied optimal spectrum allocation, in
cellular networks limited to two tiers of BSs, without flexible user
association.
Assuming a fixed number of UEs and BSs and without considering their random
spatial patterns, [15, 16, 17, 18, 19] studied optimal user association with
deterministic utility optimization. With a stochastic geometric approach, [2]
proposed the flexible user association model with bias values, which is
adopted in our work. It also derived the coverage probability and UE data
rate, considering cross-tier interference, but in non-closed forms. [2] did
not provide any design insights to optimize the derived performance metrics.
It was later extended in [3] to study optimal user association in a network
with two tiers of BSs, without considering spectrum allocation.
Joint spectrum allocation and user association was studied in [20]. It was
limited to two tiers of BSs. [20] presented a qualitative study on the optimal
network performance in terms of coverage probability and data rate, without
further providing analytical details in solving the optimization problems. A
similar problem was also studied in [21], with frequency reuse instead of
tiered spectrum division as the approach for spectrum sharing. It could
accommodate more than two tiers, but it provided only conditionally optimal
user association given frequency reuse factors or conditionally optimal
frequency reuse factors given user association bias values. Joint optimization
remained an open problem.
Compared with the above studies, we consider multiple tiers of BSs with
disjoint spectrum and provide optimal and analytically bounded near-optimal
solutions for joint spectrum allocation and user association.
## III System Model
### III-A Multi-tier Cellular Network
We consider a heterogeneous cellular network with randomly spatially
distributed $K\geq 2$ tiers of BSs. As in conventional stochastic geometric
modeling of multi-tier cellular networks [2, 20, 3, 12, 21], each tier of BSs
independently forms a homogeneous Poisson point process (PPP) in two
dimensional Euclidean space $\mathbb{R}^{2}$. Let $\Phi_{k}$ denote the PPP
corresponding to tier-$k$ BSs, with intensity $\lambda_{k}$. Without loss of
generality, we assume that $\lambda_{1}<\lambda_{2}\ldots<\lambda_{K}$. (If
$\lambda_{i}=\lambda_{j},i\neq j$ in reality, we may approximate
$\lambda_{j}=\lambda_{i}+\xi$, where $\xi$ is arbitrarily close to $0$.) UEs
are also modeled as a homogeneous PPP $\Psi$ with intensity $\mu$, independent
of all BSs. We assume each BS is connected to the core network by separate
high-capacity wired or wireless links that have no influence on our
performance analysis. In addition, because we focus on downlink analysis, we
assume that the downlink and uplink of the system are operated in different
spectrum, so that the uplink interference and capacity have no influence on
the downlink analysis.
### III-B Power and Pathloss Model
We define the tiers of BSs by their transmission power. Let $P_{k}$ be the
transmission power of tier-$k$ BSs, which is a given parameter. If
$P_{t}(\mathbf{x})$, $P_{t}(\mathbf{x})\in\\{P_{1},P_{2},\ldots,P_{K}\\}$, is
the transmission power from a BS at $\mathbf{x}$ and $P_{r}(\mathbf{y})$ is
the received power at $\mathbf{y}$, we have
$P_{r}(\mathbf{y})=\frac{P_{t}(\mathbf{x})h_{\mathbf{x},\mathbf{y}}}{\alpha|\mathbf{x}-\mathbf{y}|^{\gamma}}$,
where $\alpha|\mathbf{x}-\mathbf{y}|^{\gamma}$ is the propagation loss
function with predetermined constants $\alpha$ and $\gamma$ (where the path
loss exponent $\gamma>2$ in practice), and $h_{\mathbf{x},\mathbf{y}}$ is the
fast fading term. Corresponding to common Rayleigh fading with power
normalization, $h_{\mathbf{x},\mathbf{y}}$ is independently exponentially
distributed with unit mean. Let $H(\cdot)$ be the cumulative distribution
function of $h_{\mathbf{x},\mathbf{y}}$.
### III-C Spectrum Allocation
In order to avoid cross-tier interference, different tiers of BSs are
allocated separated spectrum. Assume the total spectrum bandwidth is $W$. The
network operator allocates $\eta_{k}W$ to each tier-$k$ BS, where $\eta_{k}$
is the spectrum allocation factor and $\sum_{k=1}^{K}\eta_{k}=1$. Let
$\bm{\eta}=(\eta_{1},\eta_{2},\ldots,\eta_{K})$. Note that BSs in the same
tier are operated on the same spectrum.
We additionally consider the possible constraints
$\eta_{\min,k}\leq\eta_{k}\leq\eta_{\max,k}$, for $k=1,2,\ldots,K$. Clearly,
we have $\sum_{k=1}^{K}\eta_{\min,k}\leq 1\leq\sum_{k=1}^{K}\eta_{\max,k}$.
Furthermore, we assume that
$0<\eta_{\min,1}\leq\eta_{\min,2}\ldots\leq\eta_{\min,K}$ and
$0<\eta_{\max,1}\leq\eta_{\max,2}\ldots\leq\eta_{\max,K}$, i.e., the network
operator is likely (but not necessarily) to allocate more spectrum to a tier
with higher BS density. Note that this is a general condition that contains
the special case where there is no constraint on $\bm{\eta}$.
Given a specific tier-$k$ BS, it is common to assume that all its associated
UEs are equally allocated spectrum [21, 2, 6]. Hence, the per-UE assigned
spectrum bandwidth is $\beta_{k}=\eta_{k}W/N_{k}$, where $N_{k}$ is a random
variable denoting the number of UEs associated with the BS.
### III-D Coverage Probability and UE Data Rate
Let $T$ denote the minimum required Signal-to-Interference Ratio (SIR) of UEs.
The coverage probability of a UE is defined as the probability that its SIR is
no less than $T$ [7]. As in conventional wireless modeling [21, 2, 6], if a UE
experiences coverage probability $P^{\prime}$ and is allocated spectrum
bandwidth $\beta^{\prime}$, its data rate is $\beta^{\prime}\log(1+T)$ if the
SIR is no less than $T$, and its data rate is $0$ if the SIR is less than $T$
(i.e., outage occurs). Thus, the overall data rate of the UE is
$\beta^{\prime}\log(1+T)P^{\prime}$. Note that, unless otherwise stated,
$\log$ represents $\log_{2}$. Also, the thermal noise is assumed to be
negligible compared with interference.
### III-E Flexible User Association
Given that a UE is located at $\mathbf{y}$, it associates itself with the BS
that provides the maximum biased received power [2, 3, 21] as follows:
$\displaystyle\mathcal{BS}(\mathbf{y})=\arg\max_{\mathbf{x}\in\Phi_{k},\forall
k}B_{k}P_{k}|\mathbf{x}-\mathbf{y}|^{-\gamma},$ (1)
where $\mathcal{BS}(\mathbf{y})$ denotes the location of the BS associated
with the UE, and $P_{k}|\mathbf{x}-\mathbf{y}|^{-\gamma}$ is the received
power from a tier-$k$ BS located at $\mathbf{x}$, and $B_{k}$ is the
association bias, indicating the connecting preference of a UE toward tier-$k$
BSs. In this case, the resultant cell splitting forms a generalized Dirichlet
tessellation, or weighted Poisson Voronoi [22], shown in Fig. 1(b). Note that
for $B_{1},B_{2},\ldots,B_{K}$, their effects remain the same if we multiply
all of them by the same positive constant. Thus, without loss of generality,
we normalize them such that $\sum_{k=1}^{K}B_{k}=1$. Let
$\mathbf{B}=(B_{1},B_{2},\ldots,B_{K})$.
Let $A_{k}$ denote the probability that a UE associates itself with a tier-$k$
BS, and $\mathbf{A}=(A_{1},A_{2},\ldots,A_{K})$. As derived in [2], we have
$\displaystyle
A_{k}=\frac{\lambda_{k}(P_{k}B_{k})^{\frac{2}{\gamma}}}{\sum_{j=1}^{K}\lambda_{j}(P_{j}B_{j})^{\frac{2}{\gamma}}},$
(2)
and thus
$\displaystyle
B_{k}=\frac{P_{k}^{-1}(A_{k}/\lambda_{k})^{\frac{\gamma}{2}}}{\sum_{j=1}^{K}P_{j}^{-1}(A_{j}/\lambda_{j})^{\frac{\gamma}{2}}}.$
(3)
Hence, there is a one-to-one mapping between $\mathbf{A}$ and $\mathbf{B}$, so
we can view them interchangeably.
### III-F Problem Statement
We first aim to derive a closed-form expression for the average UE data rate.
Then, our objective is to maximize the average UE data rate by jointly
optimizing the spectrum allocation factors $\bm{\eta}$ and the user
association bias values $\mathbf{B}$ (or equivalently $\mathbf{A}$). Finally,
we give a pricing scheme to incentivize each UE to adopt the designed
$\mathbf{B}$.
## IV Closed-form Average UE Data Rate
In this section, we derive the average UE data rate via stochastic geometric
analysis. Consider a reference UE, termed as the typical UE, communicating
with its BS, termed as the typical BS. We are interested in the typical UE
since the average UE performance in the system is the same as the performance
of the typical UE [7]. Furthermore, due to the stationarity of UEs and BSs,
throughout this section we will re-define the coordinates so that the typical
UE is located at $\mathbf{0}$.
First, we study the coverage probability given that the typical UE is
associating with a tier-$k$ BS and their distance is $d$. In this case, the
overall interference to the typical UE is the sum interference from all
tier-$k$ BSs other than the typical BS. Let $I_{k}(d)$ denote such
interference. Then
$\displaystyle
I_{k}(d)=\sum_{\mathbf{x}\in\Phi_{k}^{\prime}}\frac{P_{k}h_{\mathbf{x},\mathbf{0}}}{\alpha|\mathbf{x}|^{\gamma}}.$
(4)
where $\Phi_{k}^{\prime}$ is the Palm point process corresponding to all
tier-$k$ BSs other than the typical BS, given that the typical BS is located
at a distance of $d$ from the typical UE. It can be shown that
$\Phi_{k}^{\prime}$ is a PPP with intensity $0$ in $\mathcal{B}(\mathbf{0},d)$
and intensity $\lambda_{k}$ in
$\mathbb{R}^{2}\backslash\mathcal{B}(\mathbf{0},d)$, where
$\mathcal{B}(\mathbf{0},d)$ denotes the disk region centered at $\mathbf{0}$
with radius $d$ [7].
The distribution of $I_{k}(d)$ is derived through its Laplace transform as
follows:
$\displaystyle\mathcal{L}_{I_{k}}(d,s)=\mathbf{E}\left[\exp\left(-\sum_{\mathbf{x}\in\Phi_{k}^{\prime}}\frac{sP_{k}h_{\mathbf{x},\mathbf{0}}}{\alpha|\mathbf{x}|^{\gamma}}\right)\right]$
$\displaystyle=$
$\displaystyle\exp\left(-\lambda_{k}\int_{\mathbb{R}^{2}\backslash\mathcal{B}(\mathbf{0},d)}\left(1-\int_{\mathbb{R}^{+}}e^{-\frac{sP_{k}h}{\alpha|\mathbf{x}|^{\gamma}}}H(\mathrm{d}h)\right)\mathrm{d}\mathbf{x}\right)$
(5) $\displaystyle=$
$\displaystyle\exp\left(-\lambda_{k}\int_{\mathbb{R}^{2}\backslash\mathcal{B}(\mathbf{0},d)}\frac{\frac{sP_{k}}{\alpha|\mathbf{x}|^{\gamma}}}{\frac{sP_{k}}{\alpha|\mathbf{x}|^{\gamma}}+1}\mathrm{d}\mathbf{x}\right)$
(6) $\displaystyle=$
$\displaystyle\exp\left(-2\pi\lambda_{k}\int_{d}^{\infty}\frac{\frac{sP_{k}r}{\alpha}}{\frac{sP_{k}}{\alpha}+r^{\gamma}}\mathrm{d}r\right),$
(7)
where (5) is obtained from the Laplace functional of PPP $\Phi_{k}^{\prime}$
[7], (6) is because $h$ is exponentially distributed with unit mean, and (7)
is through a transformation to polar coordinates.
Let $P_{cover,k}(d)$ denote the conditional coverage probability of the
typical UE (given $k$ and $d$). Then
$\displaystyle P_{cover,k}(d)=$
$\displaystyle\mathbf{P}\left(\frac{P_{k}h_{\mathbf{x}_{B},\mathbf{0}}}{\alpha
d^{\gamma}}\geq TI_{k}(d)\right)$ $\displaystyle=$
$\displaystyle\mathcal{L}_{I_{k}}(d,s)|_{s=\frac{T\alpha d^{\gamma}}{P_{k}}},$
(8)
where $\mathbf{x}_{B}$ is the coordinate of the typical BS, and
$|\mathbf{x}_{B}|=d$. Substituting (7) into (8), we have
$\displaystyle P_{cover,k}(d)=$
$\displaystyle\exp\left(-2\pi\lambda_{k}\int_{d}^{\infty}\frac{Td^{\gamma}r}{Td^{\gamma}+r^{\gamma}}\mathrm{d}r\right)$
$\displaystyle\overset{t=\frac{r^{2}}{T^{2/\gamma}d^{2}}}{=}$
$\displaystyle\exp\left(-\pi\lambda_{k}T^{\frac{2}{\gamma}}d^{2}\int_{(\frac{1}{T})^{\frac{2}{\gamma}}}^{\infty}\frac{1}{1+t^{\frac{\gamma}{2}}}\mathrm{d}t\right).$
(9)
Furthermore, the probability density function of the distance between the
typical UE and its associated tier-$k$ BS is
$\displaystyle f_{k}(d)=$
$\displaystyle\frac{2\pi\lambda_{k}}{A_{k}}d\exp\left(-\pi
d^{2}\sum_{j=1}^{K}\lambda_{j}\left(\frac{P_{j}B_{j}}{P_{k}B_{k}}\right)^{\frac{2}{\gamma}}\right)$
(10) $\displaystyle=$
$\displaystyle\frac{2\pi\lambda_{k}}{A_{k}}d\exp\left(-\pi
d^{2}\frac{\lambda_{k}}{A_{k}}\right),$ (11)
where (10) is derived in [2], and (11) is by substituting (2) into (10).
Hence, the coverage probability $P_{cover,k}$ of the typical UE associated
with a tier-$k$ BS can be computed as
$\displaystyle P_{cover,k}=$
$\displaystyle\int_{0}^{\infty}f_{k}(d)P_{cover,k}(d)\mathrm{d}d$
$\displaystyle=$
$\displaystyle\int_{0}^{\infty}\frac{2\pi\lambda_{k}}{A_{k}}d\exp\left(-\pi
d^{2}\frac{\lambda_{k}}{A_{k}}\right)$
$\displaystyle\qquad\exp\left(-\pi\lambda_{k}\left(T\right)^{\frac{2}{\gamma}}d^{2}\int_{\left(\frac{1}{T}\right)^{\frac{2}{\gamma}}}^{\infty}\frac{1}{1+t^{\gamma/2}}\mathrm{d}t\right)\mathrm{d}d$
$\displaystyle=$ $\displaystyle\frac{1}{A_{k}}\frac{1}{\frac{1}{A_{k}}+C},$
(12)
where
$C=\left(T\right)^{\frac{2}{\gamma}}\int_{\left(\frac{1}{T}\right)^{\frac{2}{\gamma}}}^{\infty}\frac{1}{1+t^{\gamma/2}}\mathrm{d}t$
is a system-level constant only related to $\gamma$ and $T$.
Let $\mathbf{E}_{0}(\beta_{k})$ denote the expected spectrum bandwidth
allocated to the typical UE (connecting to a tier-$k$ BS). Following the model
in Section III-C, $\mathbf{E}_{0}(\beta_{k})$ equals the spectrum bandwidth
allocated to the typical tier-$k$ BS divided by the average number of UEs
associated with it conditioned on the typical UE, which is
$A_{k}\mu/\lambda_{k}+1$. Hence,
$\displaystyle\mathbf{E}_{0}(\beta_{k})=\frac{\eta_{k}W}{A_{k}\mu/\lambda_{k}+1}.$
(13)
Then, by Section III-D, the conditional expected data rate of the typical UE,
given it is associated with a tier-$k$ BS, can be computed as [2, 21]111By
doing so, we slightly underestimate the average data rate because the coverage
event and $\beta_{k}$ are not completely independent. Although some efforts
have been made to approximate their correlation [3, 23], all of them are
inexact but result in tremendous mathematical complexity. In Section VII, we
show that the resultant analysis is close to actual performance via
simulations.
$\displaystyle\overline{R}_{k}=\mathbf{E}_{0}(\beta_{k})\log(1+T)P_{cover,k}.$
(14)
Finally, the average data rate of the typical UE, and hence the average data
rate per UE in the system, is
$\displaystyle\mathbf{F}=$
$\displaystyle\sum_{k=1}^{K}A_{k}\overline{R}_{k}=\sum_{k=1}^{K}A_{k}\mathbf{E}_{0}(\beta_{k})\log(1+T)P_{cover,k}$
$\displaystyle=$
$\displaystyle\sum_{k=1}^{K}\frac{\eta_{k}W\log(1+T)}{(A_{k}\mu/\lambda_{k}+1)(\frac{1}{A_{k}}+C)}.$
(15)
Note that stochastic geometric analysis often leads to non-closed forms
requiring numerical integrations (e.g., [11, 12, 2, 3]), due to the integral
form of the Laplace functional or generating functional of PPPs applied in
analysis [7, 10]. Fortunately, our derived closed-form expression for the
average UE data rate facilitates the tractability of the resultant
optimization problem.
## V Joint Optimization Problem and SSAUA
We aim to maximize the average UE data rate $\mathbf{F}$ with respect to
$\bm{\eta}$ and $\mathbf{B}$. As there is a one-to-one mapping between
$\mathbf{A}$ and $\mathbf{B}$, we study the optimization problem over
$(\bm{\eta},\mathbf{A})$ instead for analytical convenience. This is formally
stated as optimization problem $\mathbf{P}$ as follows:
$\displaystyle\underset{\bm{\eta},\mathbf{A}}{\text{maximize}}$
$\displaystyle\mathbf{F}(\bm{\eta},\mathbf{A})=\sum_{k=1}^{K}\eta_{k}M_{k}(A_{k})$
subject to
$\displaystyle\sum_{k=1}^{K}\eta_{k}=1,\quad\eta_{\min,k}\leq\eta_{k}\leq\eta_{\max,k},\forall
k,$ $\displaystyle\sum_{k=1}^{K}A_{k}=1,\quad A_{k}\geq 0,\forall k,$ (16)
where $M_{k}(A_{k})$ is defined as
$\displaystyle
M_{k}(A_{k})=\frac{1}{\left(A_{k}\mu/\lambda_{k}+1\right)\left(\frac{1}{A_{k}}+C\right)}.$
(17)
Problem $\mathbf{P}$ is non-convex and cannot be solved through a standard
method. Instead, we investigate into two important structures of the optimal
solution, termed density thresholding and priority ordering, based on which we
propose a computationally efficient Structured Spectrum Allocation and User
Association (SSAUA) approach to solve the problem.
### V-A Density Thresholding Structure
First, we define an important parameter
$\displaystyle a_{k}\triangleq\sqrt{\lambda_{k}/(\mu C)}.$ (18)
Note that $M_{k}(A_{k})$ is increasing on $[0,a_{k}]$ and decreasing on
$[a_{k},\infty)$. We further observe several useful properties of
$M_{k}(A_{k})$, which are presented in Appendix -A. Based on these properties,
we obtain the following lemma, whose proof is given in Appendix -B.
###### Lemma 1
Consider a potential solution $(\bm{\eta}^{**},\mathbf{A}^{**})$ to Problem
$\mathbf{P}$. If $\exists i\neq j$, such that $A^{**}_{i}<a_{i}$ and
$A^{**}_{j}>a_{j}$, then $(\bm{\eta}^{**},\mathbf{A}^{**})$ is not an optimal
solution.
Lemma 1 suggests that, in an optimal solution, every $A_{k}$ must be on the
same side of $a_{k}$. This directly leads to the following theorem, which is
fundamental to our optimization solution.
###### Theorem 1
(Density Thresholding) Let $(\bm{\eta}^{*},\mathbf{A}^{*})$ be an optimal
solution to Problem $\mathbf{P}$. If $\sum_{k=1}^{K}a_{k}>1$, then $\forall
k,A^{*}_{k}\leq a_{k}$; if $\sum_{k=1}^{K}a_{k}<1$, then $\forall
k,A^{*}_{k}\geq a_{k}$; if $\sum_{k=1}^{K}a_{k}=1$, then $\forall
k,A^{*}_{k}=a_{k}$.
###### Proof:
If $\sum_{k=1}^{K}a_{k}>1$, because $\sum_{k=1}^{K}A^{*}_{k}=1$, $\exists l$
such that $A^{*}_{l}<a_{l}$. This leads to $A^{*}_{k}\leq a_{k}$, $\forall k$,
according to Lemma 1. The cases where $\sum_{k=1}^{K}a_{k}<1$ and
$\sum_{k=1}^{K}a_{k}=1$ are similar. ∎
Note that, the condition $\sum_{i=1}^{K}a_{i}>1$ is equivalent to
$\sqrt{\frac{1}{C}}\left(\sum_{i=1}^{K}\sqrt{\lambda_{i}}\right)>\sqrt{\mu}$,
implying the density of UEs is sparse (compared with that of BSs). Thus, we
refer to the case $\sum_{i=1}^{K}a_{i}>1$ as the _sparse UE scenario_. On the
other hand, we refer to the case $\sum_{i=1}^{K}a_{i}<1$, which is equivalent
to
$\sqrt{\frac{1}{C}}\left(\sum_{i=1}^{K}\sqrt{\lambda_{i}}\right)<\sqrt{\mu}$,
as the _dense UE scenario_. If $\sum_{i=1}^{K}a_{i}=1$, Problem $\mathbf{P}$
can be trivially solved and is ignored in the rest of our discussion. Note
that because $a_{k}$ can be computed directly from the given parameters, one
can judge which scenario Problem $\mathbf{P}$ falls within before solving the
problem. Next, the solution to $\mathbf{P}$ will be investigated separately in
the sparse UE and dense UE scenarios.
### V-B SSAUA in the Sparse UE Scenario
In this case, the original Problem $\mathbf{P}$ becomes Problem $\mathbf{P1}$
as follows:
$\displaystyle\underset{\bm{\eta},\mathbf{A}}{\text{maximize}}$
$\displaystyle\mathbf{F}(\bm{\eta},\mathbf{A})=\sum_{k=1}^{K}\eta_{k}M_{k}(A_{k})$
subject to
$\displaystyle\sum_{k=1}^{K}\eta_{k}=1,\quad\eta_{\min,k}\leq\eta_{k}\leq\eta_{\max,k},\forall
k,$ $\displaystyle\sum_{k=1}^{K}A_{k}=1,\quad 0\leq A_{k}\leq a_{k},\forall
k.$ (19)
We first observe an important ordering property of the optimal solution to
$\mathbf{P1}$, as shown in the following lemma, whose proof is given in
Appendix -C.
###### Lemma 2
(Ordering Property) Let $\mathbf{A}^{*}$ be optimal for $\mathbf{P1}$, then
$M_{1}(A^{*}_{1})\leq M_{2}(A^{*}_{2})\leq\ldots\leq M_{K}(A^{*}_{K})$.
Next, by sequentially computing $\bm{\eta}^{*}$ as follows:
$\displaystyle\left\\{\begin{array}[]{l}\eta^{*}_{K}=\min(1-\sum_{k=1}^{K-1}\eta_{\min,k},\eta_{\max,K}),\\\
\eta^{*}_{K-1}=\min(1-\eta^{*}_{K}-\sum_{k=1}^{K-2}\eta_{\min,k},\eta_{\max,K-1}),\\\
\ldots,\\\
\eta^{*}_{l}=\min(1-\sum_{k=l+1}^{K}\eta^{*}_{k}-\sum_{k=1}^{l-1}\eta_{\min,k},\eta_{\max,l}),\\\
\ldots,\\\ \eta^{*}_{1}=\min(1-\sum_{k=2}^{K}\eta^{*}_{k},\eta_{\max,1}),\\\
\end{array}\right.$ (26)
we have the following theorem:
###### Theorem 2
(Priority Ordering) Let $\mathbf{A}^{*}$ be optimal for Problem $\mathbf{P1}$,
then $(\bm{\eta}^{*},\mathbf{A}^{*})$, where $\bm{\eta}^{*}$ is computed in
(26), is an optimal solution to $\mathbf{P1}$.
###### Proof:
Consider Problem $\mathbf{P1A}$ as follows:
$\displaystyle\underset{\bm{\eta}}{\text{maximize}}$
$\displaystyle\sum_{k=1}^{K}\eta_{k}M_{k}(A^{*}_{k})$ subject to
$\displaystyle\sum_{k=1}^{K}\eta_{k}=1,\quad\eta_{\min,k}\leq\eta_{k}\leq\eta_{\max,k},\forall
k.$ (27)
Then $\mathbf{P1A}$ is a simple linear programming problem with ordered linear
coefficients in the objective, since $M_{1}(A^{*}_{1})\leq
M_{2}(A^{*}_{2})\leq\ldots\leq M_{K}(A^{*}_{K})$ due to Lemma 2. Note that
$\bm{\eta}^{*}$ does not depend on the exact values of $\mathbf{A}^{*}$; it
only requires the ordering property as shown in Lemma 2. Also, $\bm{\eta}^{*}$
is in the feasible region due to $\sum_{k=1}^{K}\eta_{\min,k}\leq
1\leq\sum_{k=1}^{K}\eta_{\max,k}$. It is easy to verify that,
$(\bm{\eta}^{*},\mathbf{A}^{*})$ is an optimal solution to $\mathbf{P1}$. ∎
Equation (26) indicates the priority ordering structure in spectrum
allocation. We see that tier-$K$ has the highest priority in spectrum
allocation, followed by tier-$(K-1)$, and so forth.
Theorem 2 provides a means to derive an optimal $\bm{\eta}^{*}$ regardless of
the $\mathbf{A}^{*}$ values. We need one further step to derive the
corresponding optimal $\mathbf{A}^{*}$ by solving the following Problem
$\mathbf{P1B}$:
$\displaystyle\underset{\mathbf{A}}{\text{maximize}}$
$\displaystyle\sum_{k=1}^{K}\eta^{*}_{k}M_{k}(A_{k})$ subject to
$\displaystyle\sum_{k=1}^{K}A_{k}=1,\quad 0\leq A_{k}\leq a_{k},\forall k.$
(28)
Note that $\mathbf{P1B}$ is a convex programming problem, since $M_{k}(A_{k})$
is concave on $[0,a_{k}]$. Thus, $\mathbf{A}^{*}$ can be computed by a
computationally efficient algorithm, such as the interior point method. Hence
the both steps to compute the jointly optimal solution
$(\bm{\eta}^{*},\mathbf{A}^{*})$ have low computational complexity.
### V-C SSAUA in the Dense UE Scenario with Performance Bound
In this case, the original Problem $\mathbf{P}$ becomes Problem $\mathbf{P2}$
as follows:
$\displaystyle\underset{\bm{\eta},\mathbf{A}}{\text{maximize}}$
$\displaystyle\mathbf{F}(\bm{\eta},\mathbf{A})=\sum_{k=1}^{K}\eta_{k}M_{k}(A_{k})$
subject to
$\displaystyle\sum_{k=1}^{K}\eta_{k}=1,\quad\eta_{\min}\leq\eta_{k}\leq\eta_{\max},\forall
k,$ $\displaystyle\sum_{k=1}^{K}A_{k}=1,\quad A_{k}\geq a_{k},\forall k.$ (29)
Problem $\mathbf{P2}$ is more complicated compared with Problem $\mathbf{P1}$,
as $M_{k}(A_{k})$ is not concave, but an S-shaped function, in the feasible
region. Hence, $\mathbf{P2}$ generally incurs high computational complexity
even if an optimal $\bm{\eta}^{*}$ is given [24, 25].
Therefore, instead of directly solving $\mathbf{P2}$, we first approximate
$M_{k}(A_{k})$ by $\widetilde{M}_{k}(A_{k})$ defined as follows:
$\displaystyle\widetilde{M}_{k}(A_{k})=\frac{1}{\left(A_{k}\mu/\lambda_{k}\right)\left(\frac{1}{A_{k}}+C\right)}.$
(30)
Note that this approximation is reasonable because $A_{k}\mu/\lambda_{k}$ is
much larger than $1$ when $\mu$ is large (i.e., the dense UE scenario). This
observation is also supported by the small performance gap as derived in
Section V-C2. Some useful properties of $\widetilde{M}_{k}(A_{k})$ are shown
in Appendix -D.
The approximated problem is referred to as Problem $\mathbf{P2A}$, where we
simply replace the objective function of $\mathbf{P2}$ by the following:
$\displaystyle\mathbf{F^{\prime}}(\bm{\eta},\mathbf{A})=\sum_{k=1}^{K}\eta_{k}\widetilde{M}_{k}(A_{k}).$
(31)
#### V-C1 Solution to $\mathbf{P2A}$
The important ordering property still holds for Problem $\mathbf{P2A}$, as
formalized in the following lemma, whose proof is given in Appendix -E.
###### Lemma 3
(Ordering Property) Let
$(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*})$ be an optimal
solution to $\mathbf{P2A}$, then
$\widetilde{M}_{1}(\widetilde{A}^{*}_{1})\leq\widetilde{M}_{2}(\widetilde{A}^{*}_{2})\leq\ldots\leq\widetilde{M}_{K}(\widetilde{A}^{*}_{K})$.
We observe that with the same ordering property, (26) can again be adopted as
an optimal solution to $\mathbf{P2A}$ in the dense UE scenario, leading to the
following theorem:
###### Theorem 3
(Priority Ordering) Let $\widetilde{\mathbf{A}}^{*}$ be optimal for Problem
$\mathbf{P2A}$, then $(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*})$,
where $\widetilde{\bm{\eta}}^{*}$ is computed the same way as $\bm{\eta}^{*}$
in (26), is an optimal solution to $\mathbf{P2A}$.
###### Proof:
The proof is similar to that of Theorem 2. ∎
Given an optimal $\widetilde{\bm{\eta}}^{*}$ for $\mathbf{P2A}$, we find the
corresponding optimal $\widetilde{\mathbf{A}}^{*}$ for $\mathbf{P2A}$ by
solving the following Problem $\mathbf{P2B}$:
$\displaystyle\underset{\mathbf{A}}{\text{maximize}}$
$\displaystyle\sum_{k=1}^{K}\widetilde{\eta}^{*}_{k}\widetilde{M}_{k}(A_{k})$
subject to $\displaystyle\sum_{k=1}^{K}A_{k}=1,\quad A_{k}\geq a_{k},\forall
k.$ (32)
Unlike in the sparse UE scenario, here we have an explicit solution, as stated
in the following theorem:
###### Theorem 4
Given an optimal $\widetilde{\bm{\eta}}^{*}$ for $\mathbf{P2A}$ (computed the
same way as $\bm{\eta}^{*}$ in (26)), the corresponding optimal
$\widetilde{\mathbf{A}}^{*}$ can be expressed as follows:
$\displaystyle\left\\{\begin{array}[]{ll}\widetilde{A}^{*}_{k}=a_{k},&k\geq
2\\\ \widetilde{A}^{*}_{1}=1-\sum_{l=2}^{K}\widetilde{A}^{*}_{l}.\\\
\end{array}\right.$ (35)
###### Proof:
See Appendix -F. ∎
Note that both (26) and (35) can be computed with low computational
complexity.
#### V-C2 Bounding the Performance Gap
Since $(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*})$ is optimal for
$\mathbf{P2A}$ rather than $\mathbf{P2}$, we next quantify the performance gap
between $(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*})$ and an
optimal solution $(\bm{\eta}^{*},\mathbf{A}^{*})$ to $\mathbf{P2}$.
The performance gap is defined as
$\displaystyle
E=\mathbf{F}(\bm{\eta}^{*},\mathbf{A}^{*})-\mathbf{F}(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*}).$
(36)
Because
$\mathbf{F}(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*})\leq\mathbf{F}(\bm{\eta}^{*},\mathbf{A}^{*})\leq\mathbf{F}^{\prime}(\bm{\eta}^{*},\mathbf{A}^{*})\leq\mathbf{F}^{\prime}(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*})$,
we have
$\displaystyle
E\leq\mathbf{F}^{\prime}(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*})-\mathbf{F}(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*})\triangleq
E^{\prime}.$ (37)
Substituting $\widetilde{\bm{\eta}}^{*}$ and $\widetilde{\mathbf{A}}^{*}$ into
$E^{\prime}$, we have
$\displaystyle E^{\prime}=$
$\displaystyle\sum_{k=1}^{K}\frac{\widetilde{\eta}^{*}_{k}}{\frac{\widetilde{A}^{*}_{k}\mu}{\lambda_{k}}\left(\frac{\widetilde{A}^{*}_{k}\mu}{\lambda_{k}}+1\right)\left(\frac{1}{\widetilde{A}^{*}_{k}}+C\right)}.$
(38)
Therefore, the relative performance gap is bounded by
$\displaystyle\epsilon\triangleq$
$\displaystyle\frac{E}{\mathbf{F}(\bm{\eta}^{*},\mathbf{A}^{*})}\leq\frac{E^{\prime}}{\mathbf{F}(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{A}}^{*})}$
(39) $\displaystyle=$
$\displaystyle\frac{\sum_{k=1}^{K}\frac{\widetilde{\eta}^{*}_{k}}{\frac{\widetilde{A}^{*}_{k}\mu}{\lambda_{k}}\left(\frac{\widetilde{A}^{*}_{k}\mu}{\lambda_{k}}+1\right)\left(\frac{1}{\widetilde{A}^{*}_{k}}+C\right)}}{\sum_{k=1}^{K}\frac{\widetilde{\eta}^{*}_{k}}{\left(\frac{\widetilde{A}^{*}_{k}\mu}{\lambda_{k}}+1\right)\left(\frac{1}{\widetilde{A}^{*}_{k}}+C\right)}}$
(40) $\displaystyle\stackrel{{\scriptstyle(a)}}{{\leq}}$
$\displaystyle\max_{k}\frac{\lambda_{k}}{\widetilde{A}^{*}_{k}\mu}\leq\max_{k}\frac{\lambda_{k}}{a_{k}\mu}$
(41) $\displaystyle=$ $\displaystyle\sqrt{\lambda_{K}C/\mu},$ (42)
where inequality (a) is obtained by observing the common factor in the
summations in the numerator and denominator of (40). This implies that
$\epsilon$ scales as $O\left(\sqrt{\lambda_{K}/\mu}\right)$. Note that because
we are considering the dense UE scenario (i.e.,
$\sqrt{\frac{1}{C}}\left(\sum_{i=1}^{K}\sqrt{\lambda_{i}}\right)<\sqrt{\mu}$),
$O\left(\sqrt{\lambda_{K}/\mu}\right)$ is small by definition.
### V-D Estimation for Complexity of Exhaustive Search
In this subsection, we briefly discuss the complexity of the exhaustive search
approach to solve Problem $\mathbf{P}$. First, as explained in details in our
technical report [26], we observe that at least one of the optimal solutions
to $\mathbf{P}$, $(\bm{\eta}^{*},\mathbf{A}^{*})$, has the following property:
there is at most one $k\in\\{1,2,\ldots K\\}$ such that
$\eta_{\min,k}<\eta^{*}_{k}<\eta_{\max,k}$; $\forall j\neq k$, either
$\eta^{*}_{j}=\eta_{\min,j}$ or $\eta^{*}_{j}=\eta_{\max,j}$ (i.e., at the
boundary). Thus, the search for $\bm{\eta}^{*}$ needs to be performed only at
these boundary cases, leading to a complexity of $\Omega(2^{K})$. Furthermore,
in the dense UE scenario, a numerical search over all locally optimal
$\mathbf{A}$ (at least $2^{K}$ of them) is required, leading to another fold
of $\Omega(2^{K})$, i.e., $\Omega(4^{K})$ overall, in complexity. Numerical
studies on the computational complexity will be presented in Section VII.
## VI Nash Equilibrium for SSAUA
Individual UEs may behave selfishly to derive unfair advantage despite our
design of $\mathbf{B}^{*}$ (or equivalently $\mathbf{A}^{*}$). Thus, in this
section, we propose a Surcharge Pricing Scheme (SPS), such that the designed
$\mathbf{B}^{*}$ is the natural outcome of a Nash equilibrium. Note that the
designed spectrum allocation factors $\bm{\eta}^{*}$ can be maintained by the
network operator and is beyond our concern.
We consider a reference individual UE, whose association bias values are
$\mathbf{B}^{\prime}=(B_{1}^{\prime},B_{2}^{\prime},\ldots,B_{K}^{\prime})$.
Let
$\mathbf{A}^{\prime}=(A_{1}^{\prime},A_{2}^{\prime},\ldots,A_{K}^{\prime})$ be
its corresponding association probabilities. For the other UEs, suppose they
all obey the association bias values $\mathbf{B}^{*}$ assigned by the network
operator. Similar to the discussions in Section III and IV, the average data
rate of the reference UE is
$\displaystyle\mathbb{F}=\sum_{k=1}^{K}\frac{\eta^{*}_{k}W\log(1+T)}{(A^{*}_{k}\mu/\lambda_{k}+1)(\frac{1}{A_{k}^{\prime}}+C)}.$
(43)
If the reference UE performs an optimization on $\mathbb{F}$ with respect to
$\mathbf{A}^{\prime}$, the resultant optimal
$\mathbf{A}^{\prime*}=(A^{\prime*}_{1},A^{\prime*}_{2},\ldots,A^{\prime*}_{K})$
is unlikely to be the same as $\mathbf{A}^{*}$. Therefore, we add the
following Surcharge Pricing Scheme: the network operator applies a surcharge
$c_{k}$ to each UE associated with a tier-$k$ BS. Let
$\mathbf{c}=(c_{1},c_{2},\ldots,c_{K})$. In this case, the average surcharge
for the reference UE is $\sum_{k=1}^{K}c_{k}A_{k}^{\prime}$. Accordingly, the
reference UE will perform the following optimization Problem $\mathbf{P3}$:
$\displaystyle\underset{\mathbf{A}^{\prime}}{\text{maximize}}$
$\displaystyle\mathbb{F}^{\prime}=\sum_{k=1}^{K}\left(\frac{\eta^{*}_{k}W\log(1+T)}{\left(A^{*}_{k}\mu/\lambda_{k}+1\right)\left(\frac{1}{A_{k}^{\prime}}+C\right)}-c_{k}A_{k}^{\prime}\right)$
subject to $\displaystyle\sum_{k=1}^{K}A_{k}^{\prime}=1,\quad
A_{k}^{\prime}\geq 0.$ (44)
Different from $\mathbf{P}$, it can be shown that $\mathbf{P3}$ is a standard
convex optimization problem. By the KKT conditions, its optimal solution
$\mathbf{A}^{\prime*}$ satisfies
$\displaystyle\frac{H_{k}}{(1+CA^{\prime*}_{k})^{2}}-c_{k}-\nu+\theta_{k}=0,$
(45) $\displaystyle\theta_{k}A^{\prime*}_{k}=0,\quad\theta_{k}\geq 0,$ (46)
where $H_{k}=\frac{\eta^{*}_{k}W\log(1+T)}{A^{*}_{k}\mu/\lambda_{k}+1}$,
$\theta_{k}$ is a Lagrange multiplier corresponding to the inequality
constraint $A_{k}^{\prime}\geq 0$, and $\nu$ is a Lagrange multiplier
corresponding to the equality constraint $\sum_{k=1}^{K}A_{k}^{\prime}=1$.
Setting $A^{\prime*}_{k}=A^{*}_{k}$, we have
$\displaystyle c_{k}=\begin{cases}\infty,&\textrm{if }A^{*}_{k}=0,\\\
\frac{H_{k}}{(1+CA^{*}_{k})^{2}}-\nu,&\textrm{otherwise}.\end{cases}$ (47)
Note that $\nu$ could be set arbitrarily due to the equality constraint.
Without loss of generality, we set
$\nu=\min_{k}\frac{H_{k}}{(1+CA^{*}_{k})^{2}}$ so that the minimum surcharge
among tiers is $0$. As a consequence, a Nash Equilibrium is achieved where
every UE adopts the assigned $\mathbf{B}^{*}$.
## VII Numerical Study
Figure 2: Average UE data rate under different UE density $\mu$.
(a) Sparse UE scenario. (b) Dense UE scenario.
Figure 3: Average UE data rate under different path loss exponent $\gamma$.
Figure 4: Designed association bias values under different UE density $\mu$.
Figure 5: Pricing values under different UE density $\mu$. Figure 6:
Comparison of run time.
In this section, we present numerical studies on the performance of SSAUA. We
label the SSAUA solution as
$(\widehat{\bm{\eta}}^{*},\widehat{\mathbf{B}}^{*})$ and
$(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{B}}^{*})$ in the sparse and
dense UE scenarios, respectively. Note that
$(\widehat{\bm{\eta}}^{*},\widehat{\mathbf{B}}^{*})$ is optimal in the sparse
UE scenario. We also compare
$(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{B}}^{*})$ with an optimal
solution $(\bm{\eta}^{*},\mathbf{B}^{*})$ obtained from exhaustive search in
the dense UE scenario.
First, we study the network performance under different UE density $\mu$. The
network parameters are as follows: $K=3$, $\lambda_{1}=1$ units/km2,
$\lambda_{2}=5$ units/km2, $\lambda_{3}=10$ units/km2, $P_{1}=56$ dBm,
$P_{2}=46$ dBm, $P_{3}=36$ dBm, $\eta_{\min,1}=0.2$, $\eta_{\min,2}=0.25$,
$\eta_{\min,3}=0.3$, $\eta_{\max,1}=0.35$, $\eta_{\max,2}=0.4$,
$\eta_{\max,3}=0.45$, $\gamma=4$, $W=200$ MHz, and $T=0.2$. In each round of
simulation, UEs and BSs are generated on a $10$ km $\times$ $10$ km square,
and the UEs in the central $5$ km $\times$ $5$ km square are sampled for
performance evaluation (in order to remove the edge effect). Each simulation
data point is averaged over all sampled UEs during 100 rounds of simulations.
The results are shown in Fig. 2. A vertical line indicates the threshold value
of $\mu$, as given in Theorem 1, separating the sparse and dense UE scenarios.
For both scenarios, we show results of the analytical and simulated
performance of SSAUA, and analytical and simulated performance of a “max
power” approach, which employs equal spectrum allocation and user association
based on the maximum received power. Since SSAUA is not optimal in the dense
UE scenario, we also add two sets of results accordingly: the analytical
optimal performance $\mathbf{F}(\bm{\eta}^{*},\mathbf{B}^{*})$ through
exhaustive search and its analytical upper bound
$\mathbf{F^{\prime}}(\widetilde{\bm{\eta}}^{*},\widetilde{\mathbf{B}}^{*})$.
Fig. 2 illustrates that the SSAUA achieves near-optimal solution in the dense
UE scenario. Furthermore, SSAUA substantially outperforms the max power
solution. Finally, the analytical performance is only slightly smaller than
the simulated performance, matching our discussions in Section IV.
Fig. 3 shows the optimal network performance under different path loss
exponent $\gamma$. The network parameters are the same as those used in Fig. 2
except $\mu$ is fixed at $100$ (i.e., sparse UE) and $500$ (i.e., dense UE)
units/km2 in Fig. 3 (a) and (b) respectively. This figure further confirms the
observations from Fig. 2. Furthermore, it shows that SSAUA is effective for a
wide range of path loss conditions.
Fig. 4 shows $\widehat{\mathbf{B}}^{*}$, $\widetilde{\mathbf{B}}^{*}$, and
$\mathbf{B}^{*}$; and Fig. 5 shows their corresponding prices
$\widehat{\mathbf{c}}^{*}$, $\widetilde{\mathbf{c}}^{*}$, and
$\mathbf{c}^{*}$, under different $\mu$. We observe that
$\widetilde{\mathbf{B}}^{*}$ and $\widetilde{\mathbf{c}}^{*}$ computed based
on SSAUA approach are close to their counterparts $\mathbf{B}^{*}$ and
$\mathbf{c}^{*}$.
Finally, a run time experiment is conducted to compare the computational
complexity of SSAUA with that of exhaustive search. The experiment is executed
by Matlab R2011a on an ASUS PC with Intel i7-3610QM 2.3GHz processor and 4GB
RAM. The results are averaged over 1000 runs for SSAUA and 10 runs for
exhaustive search (both with randomly generated parameters). Fig. 6 shows that
the run time of SSAUA is almost negligible compared with exhaustive search.
Note that the $y$-axis is in log scale. When $K$ increases, the run time of
exhaustive search exhibits an exponential growth tendency, while SSAUA remains
computationally efficient.
## VIII Conclusions
In this work, we provide a theoretical framework to study the joint
optimization of spectrum allocation and user association in heterogeneous
cellular networks. We establish a stochastic geometric model that captures the
random spatial patterns of BSs and UEs, and a closed-form expression of the
analytical average UE data rate is derived. We then consider the problem of
maximizing the average UE data rate by jointly optimizing the spectrum
allocation factors $\bm{\eta}$ and user association bias values $\mathbf{B}$,
which is non-convex programming in nature. We propose the SSAUA approach to
solve this problem with low computational complexity. We show that the SSAUA
approach is optimal in the sparse UE scenario and near-optimal in the dense UE
scenario, with a quantified tight bound scaling as
$O(\sqrt{\lambda_{K}/\mu})$. We also propose the SPS such that the designed
association bias values can be achieved in Nash equilibrium.
### -A Useful Properties of $M_{k}(A_{k})$
(M-1)
$M_{k}(A_{k})$ is increasing on $[0,a_{k}]$ and decreasing on
$[a_{k},\infty)$.
(M-2)
$M_{k}(A_{k})$ is concave on $[0,a^{\prime}_{k}]$ and convex on
$[a^{\prime}_{k},\infty)$, where $a^{\prime}_{k}$ is some threshold value
$a^{\prime}_{k}>a_{k}$. $M_{k}(A_{k})$ is concave on $[0,a_{k}]$.
(M-3)
If $\lambda_{i}<\lambda_{j}$, then $M_{i}(A)<M_{j}(A),\forall A>0$.
(M-4)
If $\lambda_{i}<\lambda_{j}$, then $M_{j}(A)-M_{i}(A)$ is strictly increasing
on $[0,a_{j}]$.
### -B Proof of Lemma 1
###### Proof:
Suppose $\mathbf{A}^{**}$ is optimal, $A^{**}_{i}<a_{i}$, and
$A^{**}_{j}>a_{j}$. Consider that we increase $A^{**}_{i}$ by a small value
$\Delta>0$ and decrease $A^{**}_{j}$ by $\Delta$. According to property (M-1),
$\eta^{**}_{i}M_{i}(A^{**}_{i})+\eta^{**}_{j}M_{j}(A^{**}_{j})<\eta^{**}_{i}M_{i}(A^{**}_{i}+\Delta)+\eta^{**}_{j}M_{j}(A^{**}_{j}-\Delta)$.
Thus, through replacing $A^{**}_{i}$ and $A^{**}_{j}$ by $A^{**}_{i}+\Delta$
and $A^{**}_{j}-\Delta$ respectively, we find a better solution to
$\mathbf{P}$, leading to a contradiction. ∎
### -C Proof of Lemma 2
###### Proof:
Suppose $\exists i<j$ such that $M_{i}(A^{*}_{i})>M_{j}(A^{*}_{j})$. This
implies that $A^{*}_{i}>A^{*}_{j}$. (Otherwise, suppose $A^{*}_{i}\leq
A^{*}_{j}$, then we have $M_{i}(A^{*}_{i})\leq
M_{i}(A^{*}_{j})<M_{j}(A^{*}_{j})$, leading to a contradiction.) A
corresponding diagram is shown in Fig. 7.
Case 1: $\eta^{*}_{i}\leq\eta^{*}_{j}$.
Let $\widehat{A}^{*}_{j}=A^{*}_{i}$ and $\widehat{A}^{*}_{i}=A^{*}_{j}$, then
we have
$\displaystyle[\eta^{*}_{i}M_{i}(\widehat{A}^{*}_{i})+\eta^{*}_{j}M_{j}(\widehat{A}^{*}_{j})]-[\eta^{*}_{i}M_{i}(A^{*}_{i})+\eta^{*}_{j}M_{j}(A^{*}_{j})]$
$\displaystyle=$
$\displaystyle[\eta^{*}_{i}M_{i}(A^{*}_{j})+\eta^{*}_{j}M_{j}(A^{*}_{i})]-[\eta^{*}_{i}M_{i}(A^{*}_{i})+\eta^{*}_{j}M_{j}(A^{*}_{j})]$
$\displaystyle=$
$\displaystyle\eta^{*}_{j}[M_{j}(A^{*}_{i})-M_{j}(A^{*}_{j})]+\eta^{*}_{i}[M_{i}(A^{*}_{j})-M_{i}(A^{*}_{i})]$
$\displaystyle\geq$
$\displaystyle\eta^{*}_{i}[M_{j}(A^{*}_{i})-M_{j}(A^{*}_{j})+M_{i}(A^{*}_{j})-M_{i}(A^{*}_{i})]>0,$
(48)
where (48) is due to property (M-4).
As a consequence, if $A^{*}_{i}$ and $A^{*}_{j}$ are replaced by
$\widehat{A}^{*}_{i}$ and $\widehat{A}^{*}_{j}$ respectively, we obtain a
larger $\mathbf{F}$, leading to a contradiction.
Case 2: $\eta^{*}_{i}>\eta^{*}_{j}$.
Let $\widehat{A}^{*}_{j}=A^{*}_{i}$, $\widehat{A}^{*}_{i}=A^{*}_{j}$,
$\widehat{\eta}^{*}_{j}=\eta^{*}_{i}$, and
$\widehat{\eta}^{*}_{i}=\eta^{*}_{j}$. (Note that because
$\eta_{\min,i}\leq\eta_{\min,j}$ and $\eta_{\max,i}\leq\eta_{\max,j}$,
$\widehat{\eta}^{*}_{j}$ and $\widehat{\eta}^{*}_{i}$ are guaranteed to be in
the feasible region.)
$\displaystyle[\widehat{\eta}^{*}_{i}M_{i}(\widehat{A}^{*}_{i})+\widehat{\eta}^{*}_{j}M_{j}(\widehat{A}^{*}_{j})]-[\eta^{*}_{i}M_{i}(A^{*}_{i})+\eta^{*}_{j}M_{j}(A^{*}_{j})]$
$\displaystyle=$
$\displaystyle[\eta^{*}_{j}M_{i}(A^{*}_{j})+\eta^{*}_{i}M_{j}(A^{*}_{i})]-[\eta^{*}_{i}M_{i}(A^{*}_{i})+\eta^{*}_{j}M_{j}(A^{*}_{j})]$
$\displaystyle=$
$\displaystyle\eta^{*}_{i}[M_{j}(A^{*}_{i})-M_{i}(A^{*}_{i})]+\eta^{*}_{j}[M_{i}(A^{*}_{j})-M_{j}(A^{*}_{j})]$
$\displaystyle>$
$\displaystyle\eta^{*}_{j}[M_{j}(A^{*}_{i})-M_{i}(A^{*}_{i})+M_{i}(A^{*}_{j})-M_{j}(A^{*}_{j})]>0.$
(49)
Thus, if $A^{*}_{i}$, $A^{*}_{j}$, $\eta^{*}_{i}$, and $\eta^{*}_{j}$ are
replaced by $\widehat{A}^{*}_{i}$, $\widehat{A}^{*}_{j}$,
$\widehat{\eta}^{*}_{i}$, and $\widehat{\eta}^{*}_{j}$ respectively, we can
find a larger $\mathbf{F}$, leading to a contradiction. ∎
Figure 7: Diagram of $M_{i}(\cdot)$ and $M_{j}(\cdot)$.
### -D Useful Properties of $\widetilde{M}_{k}(A_{k})$
(M-1’)
$\widetilde{M}_{k}(A_{k})$ is a decreasing convex function.
(M-2’)
If $\lambda_{i}<\lambda_{j}$,
$\widetilde{M}_{i}(a_{i})<\widetilde{M}_{j}(a_{j})$.
(M-3’)
If $\lambda_{i}<\lambda_{j}$, then $\widetilde{M}_{j}(A)-\widetilde{M}_{i}(A)$
is a strictly decreasing function.
(M-4’)
$\widetilde{M}_{k}(A)-\widetilde{M}_{k}(A+D)>\widetilde{M}_{k}(A^{\prime})-\widetilde{M}_{k}(A^{\prime}+D)$,
for any $A^{\prime}>A\geq a_{k}$ and $D>0$.
(M-5’)
If $\lambda_{i}<\lambda_{j}$, then
$\widetilde{M}_{j}(a_{j})-\widetilde{M}_{j}(a_{j}+D)>\widetilde{M}_{i}(a_{i})-\widetilde{M}_{i}(a_{i}+D)$,
for any $D>0$.
(M-6’)
If $\lambda_{i}<\lambda_{j}$, then
$\widetilde{M}_{j}(a_{j})-\widetilde{M}_{j}(a_{j}+D)>\widetilde{M}_{i}(A^{\prime})-\widetilde{M}_{i}(A^{\prime}+D)$,
for any $D>0$ and $A^{\prime}>a_{i}$ (combining (M-4’) and (M-5’)).
### -E Proof of Lemma 3
###### Proof:
Suppose that $\exists i<j$ (i.e., $\lambda_{i}<\lambda_{j}$) such that
$\widetilde{M}_{i}(\widetilde{A}^{*}_{i})>\widetilde{M}_{j}(\widetilde{A}^{*}_{j})$,
which also implies that
$a_{i}\leq\widetilde{A}^{*}_{i}<\widetilde{A}^{*}_{j}$. The corresponding
diagrams are shown in Fig. 8.
Case 1: $\widetilde{\eta}^{*}_{i}\leq\widetilde{\eta}^{*}_{j}$.
Case 1.1: $\widetilde{A}^{*}_{i}\geq a_{j}$.
Let $\widehat{A}^{*}_{j}=\widetilde{A}^{*}_{i}$ and
$\widehat{A}^{*}_{i}=\widetilde{A}^{*}_{j}$, then we have
$\displaystyle[\widetilde{\eta}^{*}_{i}\widetilde{M}_{i}(\widehat{A}^{*}_{i})+\widetilde{\eta}^{*}_{j}\widetilde{M}_{j}(\widehat{A}^{*}_{j})]-[\widetilde{\eta}^{*}_{i}\widetilde{M}_{i}(\widetilde{A}^{*}_{i})+\widetilde{\eta}^{*}_{j}\widetilde{M}_{j}(\widetilde{A}^{*}_{j})]$
$\displaystyle=$
$\displaystyle[\widetilde{\eta}^{*}_{i}\widetilde{M}_{i}(\widetilde{A}^{*}_{j})+\widetilde{\eta}^{*}_{j}\widetilde{M}_{j}(\widetilde{A}^{*}_{i})]-[\widetilde{\eta}^{*}_{i}\widetilde{M}_{i}(\widetilde{A}^{*}_{i})+\widetilde{\eta}^{*}_{j}\widetilde{M}_{j}(\widetilde{A}^{*}_{j})]$
$\displaystyle=$
$\displaystyle\widetilde{\eta}^{*}_{j}[\widetilde{M}_{j}(\widetilde{A}^{*}_{i})-\widetilde{M}_{j}(\widetilde{A}^{*}_{j})]+\widetilde{\eta}^{*}_{i}[\widetilde{M}_{i}(\widetilde{A}^{*}_{j})-\widetilde{M}_{i}(\widetilde{A}^{*}_{i})]$
$\displaystyle\geq$
$\displaystyle\widetilde{\eta}^{*}_{i}[\widetilde{M}_{j}(\widetilde{A}^{*}_{i})-\widetilde{M}_{j}(\widetilde{A}^{*}_{j})+\widetilde{M}_{i}(\widetilde{A}^{*}_{j})-\widetilde{M}_{i}(\widetilde{A}^{*}_{i})]>0,$
(50)
where (50) is due to property (M-3’).
Thus, if $\widetilde{A}^{*}_{i}$ and $\widetilde{A}^{*}_{j}$ are replaced by
$\widehat{A}^{*}_{i}$ and $\widehat{A}^{*}_{j}$ respectively, we obtain a
larger $\mathbf{F}^{\prime}$, leading to a contradiction.
Case 1.2: $\widetilde{A}^{*}_{i}<a_{j}$.
Let $\widehat{A}^{*}_{j}=a_{j}$, $D=\widetilde{A}^{*}_{j}-a_{j}$ and
$\widehat{A}^{*}_{i}=\widetilde{A}^{*}_{i}+D$, then we have
$\displaystyle[\widetilde{\eta}^{*}_{i}\widetilde{M}_{i}(\widehat{A}^{*}_{i})+\widetilde{\eta}^{*}_{j}\widetilde{M}_{j}(\widehat{A}^{*}_{j})]-[\widetilde{\eta}^{*}_{i}\widetilde{M}_{i}(\widetilde{A}^{*}_{i})+\widetilde{\eta}^{*}_{j}\widetilde{M}_{j}(\widetilde{A}^{*}_{j})]$
$\displaystyle=$
$\displaystyle[\widetilde{\eta}^{*}_{i}\widetilde{M}_{i}(\widetilde{A}^{*}_{i}+D)+\widetilde{\eta}^{*}_{j}\widetilde{M}_{j}(a_{j})]-[\widetilde{\eta}^{*}_{i}\widetilde{M}_{i}(\widetilde{A}^{*}_{i})+\widetilde{\eta}^{*}_{j}\widetilde{M}_{j}(\widetilde{A}^{*}_{j})]$
$\displaystyle=$
$\displaystyle\widetilde{\eta}^{*}_{j}[\widetilde{M}_{j}(a_{j})-\widetilde{M}_{j}(\widetilde{A}^{*}_{j})]+\widetilde{\eta}^{*}_{i}[\widetilde{M}_{i}(\widetilde{A}^{*}_{i}+D)-\widetilde{M}_{i}(\widetilde{A}^{*}_{i})]$
$\displaystyle\geq$
$\displaystyle\widetilde{\eta}^{*}_{i}[\widetilde{M}_{j}(a_{j})-\widetilde{M}_{j}(\widetilde{A}^{*}_{j})+\widetilde{M}_{i}(\widetilde{A}^{*}_{i}+D)-\widetilde{M}_{i}(\widetilde{A}^{*}_{i})]>0,$
(51)
where (51) is due to property (M-6’).
Thus, if $\widetilde{A}^{*}_{i}$ and $\widetilde{A}^{*}_{j}$ are replaced by
$\widehat{A}^{*}_{i}$ and $\widehat{A}^{*}_{j}$ respectively, we obtain a
larger $\mathbf{F}^{\prime}$, leading to a contradiction.
Case 2: $\widetilde{\eta}^{*}_{i}>\widetilde{\eta}^{*}_{j}$. Case 2.1:
$\widetilde{A}^{*}_{i}\geq a_{j}$. Case 2.2: $\widetilde{A}^{*}_{i}<a_{j}$.
The proof for these cases can be found in our technical report [26]. ∎
(a) Case 1.1 and 2.1.
(b) Case 1.2 and 2.2.
Figure 8: Diagrams of $\widetilde{M}_{i}(\cdot)$ and
$\widetilde{M}_{j}(\cdot)$.
### -F Proof of Theorem 4
###### Proof:
Suppose $\exists k\geq 2$ such that $\widetilde{A}^{*}_{k}>a_{k}$. Let $l=1$,
$\widehat{A}^{*}_{k}=a_{k}$, $D=\widetilde{A}^{*}_{k}-a_{k}$, and
$\widehat{A}^{*}_{l}=\widetilde{A}^{*}_{l}+D$. Similar to the proof of Lemma
3, we can show that if we replace $\widetilde{A}^{*}_{l}$ and
$\widetilde{A}^{*}_{k}$ by $\widehat{A}^{*}_{l}$ and $\widehat{A}^{*}_{k}$
respectively, we find a better solution to Problem $\mathbf{P2A}$, which leads
to a contradiction. See [26] for details.
∎
## References
* [1] A. Damnjanovic, J. Montojo, Y. Wei, T. Ji, T. Luo, M. Vajapeyam, T. Yoo, O. Song, and D. Malladi, “A survey on 3GPP heterogeneous networks,” _IEEE Wireless Communications_ , vol. 18, no. 3, pp. 10–21, Jun. 2011.
* [2] H.-S. Jo, Y. J. Sang, P. Xia, and J. G. Andrews, “Heterogeneous cellular networks with flexible cell association: A comprehensive downlink SINR analysis,” _IEEE Trans. on Wireless Communications_ , vol. 11, no. 10, pp. 3484–3495, Oct. 2012.
* [3] S. Singh, H. Dhillon, and J. G. Andrews, “Offloading in heterogeneous networks: Modeling, analysis, and design insights,” _IEEE Trans. on Wireless Communications_ , vol. 12, no. 5, pp. 2484–2497, May 2013.
* [4] W. C. Cheung, T. Q. S. Quek, and M. Kountouris, “Stochastic analysis of two-tier networks: Effect of spectrum allocation,” in _Proc. of IEEE International Conference on Acoustics, Speech and Signal Processing (ICASSP)_ , Prague, Czech Republic, May 2011.
* [5] V. Chandrasekhar and J. G. Andrews, “Spectrum allocation in tiered cellular networks,” _IEEE Trans. on Communications_ , vol. 57, no. 10, pp. 3059–3068, Oct. 2009.
* [6] W. C. Cheung, T. Q. S. Quek, and M. Kountouris, “Throughput optimization, spectrum allocation, and access control in two-tier femtocell networks,” _IEEE Journal on Selected Areas in Communications_ , vol. 30, no. 3, pp. 561–574, Apr. 2012.
* [7] F. Baccelli and B. Blaszczyszyn, “Stochastic geometry and wireless networks, volume 1: Theory,” _Foundations and Trends in Networking_ , vol. 3, no. 3-4, pp. 249 – 449, 2009.
* [8] ——, “Stochastic geometry and wireless networks, volume 2: Applications,” _Foundations and Trends in Networking_ , vol. 4, no. 1-2, pp. 1–312, 2009\.
* [9] M. Haenggi, J. G. Andrews, F. Baccelli, O. Dousse, and M. Franceschetti, “Stochastic geometry and random graphs for the analysis and design of wireless networks,” _IEEE Journal on Selected Areas in Communications_ , vol. 27, no. 7, pp. 1029 – 1046, Sept. 2009.
* [10] D. Stoyan, W. Kendall, and J. Mecke, _Stochastic Geometry and Its Applications_ , 2nd ed. Wiley, 1995.
* [11] J. G. Andrews, F. Baccelli, and R. Ganti, “A tractable approach to coverage and rate in cellular networks,” _IEEE Trans. on Communications_ , vol. 59, no. 11, pp. 3122–3134, Nov. 2011.
* [12] H. Dhillon, R. Ganti, F. Baccelli, and J. G. Andrews, “Modeling and analysis of K-tier downlink heterogeneous cellular networks,” _IEEE Journal on Selected Areas in Communications_ , vol. 30, no. 3, pp. 550–560, Apr. 2012.
* [13] H.-S. Jo, P. Xia, and J. G. Andrews, “Open, closed, and shared access femtocells in the downlink,” _arXiv:1009.3522 [cs.NI]_ , 2010.
* [14] W. C. Cheung, T. Q. S. Quek, and M. Kountouris, “Spectrum allocation and optimization in femtocell networks,” in _Proc. of IEEE ICC_ , Ottawa, Canada, June 2012.
* [15] J.-W. Lee, R. Mazumdar, and N. Shroff, “Joint resource allocation and base-station assignment for the downlink in CDMA networks,” _IEEE/ACM Trans. on Networking_ , vol. 14, no. 1, pp. 1–14, Feb. 2006.
* [16] Q. Ye, B. Rong, Y. Chen, M. Al-Shalash, C. Caramanis, and J. G. Andrews, “User association for load balancing in heterogeneous cellular networks,” _IEEE Trans. on Wireless Communications_ , vol. 12, no. 6, pp. 2706–2716, Jun. 2013.
* [17] M. Hong and A. Garcia, “Mechanism design for base station association and resource allocation in downlink OFDMA network,” _IEEE Journal on Selected Areas in Communications_ , vol. 30, no. 11, pp. 2238–2250, Dec. 2012\.
* [18] C. S. Chen, F. Baccelli, and L. Roullet, “Joint optimization of radio resources in small and macro cell networks,” in _Proc. of Vehicular Technology Conference (Spring)_ , Budapest, Hungary, May 2011.
* [19] K. Shen and W. Yu, “Downlink cell association optimization for heteregeneous networks via dual coordinate descent,” in _Proc. of IEEE International Conference on Acoustics, Speech and Signal Processing (ICASSP)_ , Vancouver, Canada, May 2013.
* [20] S. Singh and J. G. Andrews, “Joint resource partitioning and offloading in heterogeneous cellular networks,” _arXiv:1303.7039 [cs.IT]_.
* [21] Y. Lin and W. Yu, “Optimizing user association and frequency reuse for heterogeneous network under stochastic model,” _in Proc. of IEEE Globecom, Atlanta, GA, 2013_.
* [22] P. F. Ash and E. D. Bolker, “Generalized Dirichlet tessellations,” _Geometriae Dedicata_ , vol. 20, no. 2, pp. 209–243, Apr. 1986.
* [23] J.-S. Ferenc and Z. Neda, “On the size distribution of Poisson Voronoi cells,” _Physica A: Statistical Mechanics and its Applications_ , vol. 385, no. 2, pp. 518–526, Nov. 2007.
* [24] M. Udell and S. Boyd, “Maximizing a sum of sigmoids.” [Online]. Available: http://www.stanford.edu/~udell/doc/max_sum_sigmoids.pdf
* [25] M. Chiang, “Nonconvex optimization of communication systems,” _Advances in Mechanics and Mathematics, Special Volumn on Strang’s 70th Birthday_ , pp. 137–196, Oct. 2007.
* [26] W. Bao and B. Liang, “Structured spectrum allocation and user association in heterogeneous cellular networks,” University of Toronto, Tech. Rep., 2013. [Online]. Available: http://www.comm.utoronto.ca/~liang/publications/techreport/SSAUA.pdf
|
arxiv-papers
| 2013-09-29T03:48:01 |
2024-09-04T02:49:51.652120
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Wei Bao and Ben Liang",
"submitter": "Wei Bao",
"url": "https://arxiv.org/abs/1309.7527"
}
|
1309.7563
|
A closed parameterization of DNA–damage by charged particles, as a function of
energy — A geometrical approach
Frank Van den Heuvel PhD∗1,2,
1 CRUK/MRC Oxford Institute for Radiation Oncology, Department of Oncology,
University of Oxford, Oxford, UK
2 Laboratory for experimental radiotherapy, Department of Oncology, University
of Leuven, Leuven, Belgium
$\ast$ E-mail: [email protected]
## Abstract
Purpose: To present a closed formalism calculating charged particle radiation
damage induced in DNA. The formalism is valid for all types of charged
particles and due to its closed nature is suited to provide fast conversion of
dose to DNA-damage.
Methods: The induction of double strand breaks in DNA–strings residing in
irradiated cells is quantified using a single particle model. This leads to a
proposal to use the cumulative Cauchy distribution to express the mix of high
and low LET type damage probability generated by a single particle. A
microscopic phenomenological Monte Carlo code is used to fit the parameters of
the model as a function of kinetic energy related to the damage to a DNA
molecule embedded in a cell. The model is applied for four particles:
electrons, protons, alpha–particles, and carbon ions. A geometric
interpretation of this observation using the impact ionization mean free path
as a quantifier, allows extension of the model to very low energies.
Results: The mathematical expression describes the model adequately using a
chi–square test ($\chi^{2}/NDF<1$). This applies to all particle types with an
almost perfect fit for protons, while the other particles seem to result in
some discrepancies at very low energies. The implementation calculating a
strict version of the RBE based on complex damage alone is corroborated by
experimental data from the measured RBE. The geometric interpretation
generates a unique dimensionless parameter $k$ for each type of charged
particle. In addition, it predicts a distribution of DNA damage which is
different from the current models.
## Introduction
The biological effect of ionizing radiation on human cells is believed to be
related to the generation of damage in the DNA–molecule located in the cell’s
nucleus[1]. The physical mechanism is the ionization of the DNA macro
molecule, generating lesions in the molecular structure, either by direct
ionization or by the generation of radicals in the vicinity of the DNA which
then indirectly damage it. These events (direct or indirect) can create
several types of damage to the DNA by combining a number of lesions into a
cluster, which can only happen if they occur in close proximity (typically
within one turn of the DNA–helix). The most prevalent of these damage types
are base damage (2 lesions), followed by single strand breaks (SSB) (3
lesions) , double strand breaks (DSB) (4 lesions), and locally multiple damage
sites (LMDS). The latter are clusters of different types of damage occurring
close to each other. It is shown that base as well as SSB damage is not likely
to be a deciding factor in the destruction of cells, due to the efficient
repair mechanisms which exist in the cell[2]. The combination of double strand
breaks and LMDS’s is likely to be the root cause for cell kill[3].
To quantify the amount of ionizing interactions in a medium, the physical
notion of dose can be used. Dose is defined as the amount of energy deposited
in a medium per unit mass and is expressed in Joule(J) per kg or Gray (Gy). In
the case of dose deposition by charged particles the Bethe–formalism is used.
This describes ionization events in a medium in terms of energy loss of the
charged particles in inelastic collisions with the electrons of the medium,
through the notion of mass stopping power ($dE/\rho dx$). In his seminal work
already in 1930, Bethe showed that there is an intimate relationship between
stopping power on the one hand, and energy (i.e. speed), charge, and the
medium in which the interaction takes place on the other hand[4]. A further
extension taking into account the possibility of the charged particle picking
up electrons, thereby changing the stopping power was introduced by Barkas[5],
using the concept of an effective charge.
In radiation biology, linear energy transfer (LET) is used rather than
stopping power. LET is identical to stopping power with the energy delivered
to $\delta$–rays (i.e. highly energetic knock on electrons) subtracted. This
quantity is called restricted stopping power. As such, LET is a measure for
the density of ionization taking place along the track of a charged particle
through a medium. Due to its close relationship with stopping power, it
follows that there is a close relationship between LET and the kinetic energy
of the depositing particle. From observation a dearth of DSB’s and LMDS’s was
shown to be related to low LET irradiations, while an increased number of both
for the same dose is seen high in hight LET irradiations[1]. Brenner and
Ward[6] argued that DSB and LMDS damage was related to multiple interactions
by single particles, rather than the combination of single strand breaks
generated by single particles. In the field of microdosimetry, this is taken a
step further by defining the notion of lineal energy which introduces the
amount of energy deposited along lines confined in a convex geometric shape
with a given distribution of cord lengths estimating the energy deposited in
various shapes, which can be used for measurement (i.e. spheres, cylinders).
Extending this, it is natural to propose a model where distance between
ionizations along these lines plays a significant role in the generation of
DNA–damage. A full listing and treatment of these quantities can be found in
the ICRU reports 16, 19, and 36 [7, 8, 9].
To describe the damage impact of charged particles on the DNA–structure, the
science community has taken its recourse to using Monte Carlo simulations to
quantify the damage introduced[10, 11]. A more fundamental analytical approach
is currently lacking, due to the underlying complexity of the DNA molecule,
and the paucity of the available experimental data. The data which is
available is mainly provided in terms of relative biological effective dose
(RBE), a quantity combining physical, spectral, chemical, and biological
factors, all of which hamper ab–initio calculations.
Monte Carlo calculations are able to predict the induction of simple or
complex damage as well as induction of single and double strand breaks in
DNA–molecules. These findings are interpreted using the Bethe–Barkas formalism
in terms of LET and show that high LET particles indeed introduce more complex
damage.
In this paper we develop a parameterization using a simple geometrical model,
that describes the behavior as calculated by the Monte Carlo codes. We also
show that this formalism describes the current knowledge well.
## Methods and Materials
### Theory
We use the single charged particle model as proposed by Brenner and Ward,
distinguishing three types of interaction results: Low LET, high LET, and
intermediate LET mode. The specifics of each mode are explained below.
1. 1.
Low LET: A single particle is generally not able to generate lesions close
enough together to induce double strand breaks at each interaction. It is
clear that DSB’s can be generated but in a limited fashion and that we use the
word lesion in liberal fashion to indicate an interactive event which has
damage as a consequence.
2. 2.
High LET: The particle has the possibility to generate multiple lesions
irrespective of any geometrical considerations. We implicitly assume that the
double strand break damage is the result of multiple interactions by one
particle. How exactly this damage is introduced (direct or indirect) is
outside the scope of this article. An implicit assumption however is that
ionizing events need to be geometrically close to the DNA structure.
3. 3.
Intermediate: In given geometric circumstances it is possible for the charged
particle to generate DSB–damage, in a high–LET manner, depending on the angle
under which the particle hits the sensitive volume (Fig. 1).
As a surrogate to categorize the charged particle in one of the types defined
above, we use the mean path length between ionizing interactions in a medium
consistent with the atomic make up of a DNA–molecule for the type of particle
under consideration. In the remainder, we denote this with $\lambda(E)$, where
$E$ is the kinetic energy of the particle. If $\lambda(E)$ is large relative
to the sensitive volume, then the lesions on average are too far apart and
only damage types related to a few lesions can occur (i.e. SSB and base
damage). Charged particles with such energies will be part of the first
category. If on the other hand $\lambda(E)$ is small then the probability of
lesions creating more complex clusters of damage close together will be
higher. Charged particles with this property will be in the high LET category.
Finally, charged particles with intermediate distances between ionizing events
have the capability of generating DSB and LMDS damage depending on other
factors than $\lambda(E)$ alone. In this model we use the geometric direction
of the path of the charged particle as a parameter. In Figure 1 a schematic
model of this approach is shown. This implies that only a limited amount of
directions are available to contribute to the amount of complex damage in the
manner as outlined for the high–LET type interactions. This happens when for a
particle of a given energy the quantity $\lambda(E)$ is slightly larger than
the maximal distance between two DNA–damage lesions to be considered as being
in the same cluster (usually about 10 base pairs (bp)). Due to the finite
thickness of the sensitive volume it is possible to behave in a high–LET
fashion depending on the angle with which the particle’s path crosses the
volume. This occurs when the projection of the path is smaller than the
previously determined maximum.
Figure 1: A schematic model of a source of charged particles with a given
mean free path length (i.e. a given energy), which is comparable with the
diameter of the sensitive cell volume. As the angle ($\theta$) of the
particle’s path with respect to the normal to the axis of the structure
increases the chance that more than a single event will occur in the volume.
This implies that if the angle ($\theta$) is larger than the one for which the
projection of the average length between interactions equals the diameter,
more high LET events will be registered.
### Equivalence Principle
In the case of irradiation with charged particles all directions of the
particle’s paths are possible as are all rotational positions of the
DNA–structure. A particle that interacts (i.e. that creates a lesion) at the
surface of a given sensitive volume has limited possibilities to interact
again given that on average, a specific distance (which depends on the
particle energy) has to be travelled before it interacts again. The next
interaction’s position is then limited by the constraints outlined above if it
is to fall within the sensitive volume. This first interaction can happen
anywhere along the volume, but the constraints are relative to the position of
that point. This implies that we can invoke an equivalence principle and
reduce the problem to that of an isotropic point source positioned at the
surface of a sensitive volume.
### Mathematical expression of the equivalence principle
We need to calculate what fraction of the paths starting in the given point
can interact with the sensitive volume given the fact that there is a length
within which this is not likely, provided by $\lambda(E)$, and that there is a
maximal distance ($H$) that disqualifies the generated lesion to be registered
in the same cluster. We have reduced this problem to that of the distribution
of projections of a point source on a line–piece, the solution to which is
known as the Cauchy–distribution[12], and is described by the Lorenz function
$f(x)$ with $x\in\Re$ expressed as follows:
$f(x)~{}=~{}\frac{1}{1+(\frac{x}{r})^{2}}$ (1)
Figure 2: The abstracted version of Figure 1 describes the distribution of
horizontal distances at which a line segment tilted at a random angle $\theta$
cuts the x–axis. Only particles with large angles contribute to double strand
break events by combining damage generated by a single particle. The red line
indicates the “forbidden” area as (on average) this distance between the
ionization events is observed.
From Figure 2, it follows that the contribution $P$ for a given energy of
charged particles to high LET events is proportional to:
$P~{}\sim~{}2\int_{H-\lambda(E)}^{H}\frac{1}{1+(\frac{x}{r})^{2}}dx$ (2)
Performing the calculation, we obtain:
$\displaystyle P~{}$
$\displaystyle\sim~{}\frac{2}{\pi}[\tan^{-1}(\frac{x}{r})\biggm{|}_{H-\lambda(E)}^{H}]+C$
(3)
$\displaystyle\sim~{}\frac{2}{\pi}[\tan^{-1}(\frac{H}{r})-\tan^{-1}(\frac{H-\lambda(E)}{r})]+C$
(4)
$\displaystyle\sim~{}\frac{2}{\pi}[\tan^{-1}(\frac{\lambda(E)-H}{r})+\tan^{-1}(\frac{H}{r})]+C$
(5)
This implies that the amount of DSB–damage for a given dose and given energy
of the charged particle is governed by the following expression.
$F_{cd}(E)~{}=~{}(a-b)\frac{2}{\pi}[\tan^{-1}(\frac{\lambda(E)-H}{r})]+b$ (6)
The change from a low to a high LET regimen occurs over a small energy
interval. In such a small interval the average distance dependence on the
energy of the particle can be approximated with a linear function. Therefore,
we expect the energy dependence of the contribution of complex damage to
follow the same form as in Equation 6 yielding the following expression, with
H=$\lambda(E_{0})$ and r=$\lambda(\Gamma/2)$, $E_{0}$ being the energy, where
the change in DSB is maximized and $\Gamma$ a measure for the width of the
slope (i.e. the full width at half maximum in differential energy space).
$F_{cd}(E)~{}=~{}(a-b)\frac{2}{\pi}[\tan^{-1}(\frac{E-E_{0}}{\Gamma/2})]+b$
(7)
With the parameters $a$ and $b$ related to the levels 1 and 2 as outlined
above. From boundary conditions we find that at very large energies (i.e.
$E>>E_{0}$) the expression is reduced to minimal number of double interactions
($D_{min}$) which is equal to $a$. The value of $b$ is related to the maximal
number of double interactions ($D_{max}$) as follows:
$D_{max}=D_{min}+(1-\tan^{-1}(-\frac{E_{0}}{\Gamma/2}))b.$ (8)
The formalism using energy alone allows us to forego specific assumptions
regarding the dimensions of the DNA–molecule. Furthermore, it also allows us
to apply this technique to particles where the values for $\lambda(E)$ are
less well known. In addition, it allows us to test this formalism using
experimentally available data which is available as a function of energy.
### Validation using Monte Carlo Simulations
The use of microdosimetric calculations has provided important insight into
the mechanisms and effects of radiation deposition. In the past, Monte Carlo
simulations of charged particle deposition by various modalities were used to
quantify and typify the kinds of damage introduced by the different
modalities[10].
The Monte Carlo Damage Simulation code (MCDS) developed by Semenenko and
Stewart, generates spatial maps of the damaged nucleotides forming many types
of clustered DNA lesion, including single-strand breaks (SSB), double strand
breaks (DSB), and individual or clustered base damages[11]. This approach has
been shown to yield a linear relationship of the number of generated DSB’s up
to a high dosage. It follows that this parameterization also provides the
possibility to link dose to damage. In this paper, MCDS version 3.0 was used
with the parameters described below. The DNA length which was chosen to be
1Gbp (Giga base pairs) and a nucleus diameter of 5$\mathrm{\mu m}$. In the
MCDS software, the geometry of the DNA–molecule is not an explicit parameter.
Here four parameters are used: 1) the DNA–segment length $n_{seg}$, which is
an ad hoc parameter expressed as base pairs $Gy^{-1}cell^{-1}$, 2) the number
of strand breaks generated $\sigma_{Sb}$, 3) the number of base pair damages
generated $\sigma_{Bb}$ by defining $f=\sigma_{Bb}/\sigma_{Sb}$, and 4) a
parameter $N_{min}$ (bp) describing the minimal separation for damage to be
apart not to be counted as being in the same cluster. The values of these
parameters is determined on the basis of other simulations and measurements.
For a more in–depth treatment of these parameters we refer to the work by
Semenenko and Stewart[13]. Variable input parameters MCDS were; the modality
(i.e. energy depositing particle (electron, proton,…)), the energy (in MeV),
and the oxygen concentration in %. In the implementation described here we
chose to omit any oxygen enhancement as this could be a confounding factor and
is the subject of another study. In this study it was found that oxygen only
changed the amount of damage in the low LET regimen, leaving the formalism
unchanged (data not shown). Therefore, a concentration of 0% oxygen was used.
For every particle type at the relevant kinetic energies, all complex damage
was noted per Gy, per cell and per kinetic energy.
### Fitting procedure
The ultimate goal was to fit the complex damage function to the data obtained
by the Monte Carlo simulation. The parameters that need fitting are the energy
position ($E_{0}$) the width of the underlying Cauchy distribution ($\Gamma$)
and the parameters $a$ and $b$. If a regular fit (i.e. all parameters fit at
the same time) is performed we see strong co–variances between the parameters.
To come to meaningful results we opted to perform a two step procedure: First,
we eliminate the parameters $a$ and $b$ by fitting the differential, thereby
reducing expression 6 to the Lorenz function.
$\frac{dF_{cd}}{dE}=\frac{\Gamma^{2}/4}{(E-E_{0})^{2}+\Gamma^{2}/4}$ (9)
This is also mathematically equivalent to the fit of a Breit–Wigner resonance
in high energy physics[14]. In a second fit–procedure, the remaining variables
$a$ and $b$ are fit using the cumulative Cauchy function. The fitting
procedures were performed in the gnuplot111http://www.gnuplot.info–software
using a Levenberg–Marquardt minimization routine.
## Results
In Figure 3, the Lorenz expression as outlined in Equation 6 together with a
normalization factor, is used to fit the energy differential probability for
the generation of DSBs. The fit is performed to minimize the $\chi^{2}$–value.
In all cases, the resulting $\chi^{2}/NDF$ (NDF = Number of Degrees of
Freedom) are lower than 1. The values of the parameters are provided in Table
1.
(a) Electrons (b) Protons (c) $\alpha$–particles (d) Carbon Ions
Figure 3: Fitting the Cauchy expression to the energy differential probability of generating DSB’s denoted $\frac{d\sigma}{dE}$ Particle | $\Gamma$ | $E_{0}$ | $a$ | $b$
---|---|---|---|---
$e^{-}$ | (2.854$\pm$0.051)$10^{-04}$ MeV | (1.05736$\pm$0.036)$10^{-04}$ MeV | 2.9061 | 21.460
$p^{+}$ | 0.5575$\pm$0.0094 MeV | 0.1642 $\pm$ 0.0037 MeV | 2.89068 | 21.4273
$\alpha^{++}$ | 8.20$\pm$0.17 MeV | 3.1850$\pm$0.056 MeV | 3.0856 | 20.7933
C6+ | 201.7$\pm$8.4 MeV | 95.4$\pm$2.5 MeV | 3.01459 | 21.8489
Table 1: The different values for $\Gamma$ and $E_{0}$ as defined by Equation
9 and obtained from a fitting procedure together with the asymptotic standard
error of the fitted parameter. All fits exhibited minimal values of
$\chi^{2}/NDF$ (NDF = Number of Degrees of Freedom). The columns $a$ and $b$
are the parameters indicating the levels of DSB at low, resp. high LET. Note
that even in low LET the number of DSB’s is not zero as complex damage can
occur due to the combination of simple damage events.
All fits are completely satisfactory at energies higher than $E_{0}$. On the
lower energy side some discrepancies can be observed depending on the incoming
particles, particularly in the case of electrons and carbon ions. We refer the
reader to the discussion section. For protons we see a satisfactory fit over
the full energy range.
Figure 4 shows the final results with all parameters fit. Again, all fits have
$\chi^{2}$–values commensurate with a positive goodness of fit. The final
values and the standard errors for the fitted parameters are listed in Table
1. Note, that the noise in the differential curves increases as the particles
become heavier. The random-seeming errors in the estimates of the derivative
arise in part from the Monte Carlo estimates of the mean number of DSB per Gy
per Gbp and from numerical instabilities associated with the calculation of
the derivative using finite difference methods.
(a) Electrons (b) Protons (c) $\alpha$–particles (d) Carbon Ions
Figure 4: The prediction of the number of double strand breaks or more
complex damage as a function of energy for 4 relevant charged particles. This
provides the number of Double Strand breaks (DSB) per Gy, Gbp and per cell.
The prediction for protons and alpha particles is almost perfect. For
electrons and carbon ions some discrepancies exist at lower energies.
### Geometric approach
Now is the time to investigate the geometric interpretation further. To
quantify the function $\lambda(E)$ we can use the inelastic mean free path as
measure (IMFP). Values for IMFP for electrons are well known in the
literature, not in the least as they are important in solid state physics and
electron microscopy. They can be found in freely available databases for a
variety of elements and compounds, even for organic molecules like DNA[15].
Proton values can be found in a publication by Zhen–Yu and colleagues[16]. For
heavier particles such as $\alpha$–particles and carbon–ions, the data is more
difficult to find. We therefore opt not to use the data for these particles
and restrict ourselves to electrons and protons in this further treatment.
In all current microdosimetric codes, the Bethe formalism is used which is
valid for higher energies (i.e. above 500eV for electrons). This implies that
changes in IMFP, denoted by $\lambda$, which impact the damage calculated
using these codes, also reflect the limitations of the Bethe formalism. From
the theory the following expression is used:
$\lambda(E)=\frac{E}{A\log(E/E_{0})+B}$ (10)
Particle | A | B
---|---|---
Electrons | 69.200 eV/nm | -153.94 eV/nm
Protons | 115.231 keV/nm | -301.45 keV/nm
Table 2: Parameters obtained by fitting Eq. 10 to data obtained from NIST
(electrons) and Zhen Yu et al. (protons)
In this work the parameters $H$ and $r$ have thus far not been linked to any
physical property but were fit. An interesting proposition could be to link
these to dimensions of the target structure. Indeed, the choice of a cylinder
as a geometric representation is not an accident. It is natural to use the
diameter of a DNA–molecule as a measure of the cylinder’s diameter. The length
of the cylinder is then related to the maximal distance we allow to classify
two damage events, being part of the same cluster of complex damage. Both
values can readily be found in the literature and text books[17]. For the most
prevalent form of cellular DNA (B–DNA), the values are $3.4$nm (i.e. the
height of a spiral of 10 base pairs), and $2.37$nm as the diameter. We now
define a dimensionless quantity $k$ which is specific to the type of charged
particle used. It is clear that this parameter acts as a scaling parameter but
also depends on the ratio of both fixed parameters. Equation 6 now reads as
follows:
$F_{d}(E)~{}=~{}(a-b)\frac{2}{\pi}[\tan^{-1}(\frac{k\lambda(E)-3.4}{2.37})]+b$
(11)
This reduces the impact of the charged particle’s energy on the induction of
complex damage in a DNA–molecule to three parameters $a$, $b$, and $k$. Figure
5 illustrates the use of these parameters and shows that comparable results to
the energy–based formalism can be obtained. It follows that we can repeat the
fitting procedure keeping $a$ and $b$ from the expression based on energy (Eq
7). We find values of k=5.18 for electrons and k=4.82 for protons.
### Extending the model
In the work presented above as well as in the used Monte Carlo simulations,
the Bethe–Barkas formalism together with its flawed approach in the lower
energy regions has always been used. It is well established that the IMFP does
not follow the expression outlined in Equation 10, where $\lambda(E)$ keeps
diminishing as the energy diminishes. Indeed, when the energy is lower than
200eV an increase in IMFP is observed due to plasmonic effects[18]. Ziaja et
al[19] showed that it is possible to describe this behavior analytically by
extending Equation 10 with a second term as follows:
$\lambda(E)_{Z}~{}=~{}\frac{\sqrt{E}}{A_{1}(E-E_{th})^{B_{1}}}~{}+~{}\frac{E-E_{0}\exp(-B/A)}{A\log(E/E_{0})+B}$
(12)
In this equation the parameter $E_{th}$ serves as a threshold separating the
behavior as described by Bethe from the plasmonic interactions. Using the data
provided in the work from Zhen–Yu and colleagues[16] it is straightforward to
obtain parameters for the behavior of protons. These are presented in Table 3.
Particle | $A_{1}$ | $B_{1}$ | $E_{th}$ | $A$ | $B$ | $E_{0}$
---|---|---|---|---|---|---
Electrons | 0.6560 | 1.0100 | 24.2838 | 65.898 | -128.23 | 1.0
Protons | 0.681 | 1.249 | 42.38 | 117.01 | -318.7 | $1.0\times 10^{3}$
Table 3: Parameters as in Table 2 with added lower energy factors. The
fitting was performed using Eq. 12
To extend our model to incorporate the behavior of very low energy particles
it is sufficient to replace the expression $\lambda(E)$ by $\lambda(E)_{Z}$ in
equation 11. In Figure 5, the modified curves show the difference with the
calculations based on the Bethe formalism only. This also shows that there is
an upper limit to the increase in DSB’s which depends on the type of particle.
It is conceivable that this approach also works for the heavier particles
which can be seen when using the IFMP’s in water for these (not shown).
Figure 5: Using the quantities for $H$ and $r$, the dimensionless constant
$k$ for electrons (left) and protons (right) is determined. Using both the
limited expression for $\lambda(E)$ and the more accurate estimate
$\lambda_{Z}(E)$. The former provides a fit to the Monte Carlo data comparable
with the results obtained using the energy–based formalism. The second
approach provides a maximal complex damage yield which differs for electrons
and protons.
## Discussion
We developed an approach to predict damage in complicated situations where
fields of different charged particles and their respective energy spectra
impact on living cells. The approach, due to its analytical nature allows very
fast calculation of damage in otherwise long simulations. In the derivation of
this approach using energy alone there are no assumptions on the mechanics
with which DNA–damage is caused by the charged particles. The only assumption
is that there is a sensitive volume where, if ionizations take place, damage
is introduced in the DNA. How exactly this damage is caused is not specified.
In the remainder of the text a parameter is identified, the average distance
between ionizations for the given charged particle in the medium ($\lambda$).
We show that this approach adequately quantifies the results from Monte Carlo
simulations based on phenomenological data and reduces these to a closed
analytical expression whereby the type of charged particle is expressed by a
single parameter ($k$). On the other hand we should be aware that issues like
repair mechanisms and oxygen effects are not present in the model, making its
applicability limited. However, if all things are identical (i.e. the type of
cells, oxygenation, etc…) and the only thing different is the type of charged
particle and its energy, then the original damage introduced in the DNA
structure should correlate with the outcome. An underlying assumption here is
that the repair processes are somehow independent from the modality with which
the cell is irradiated.
The results of this approach can be applied to determine the biological impact
of radiation in mixed environments, as in the case of proton therapy, where
protons, electrons and heavier ions (due to neutrons), deposit energy. Other
approaches have been proposed to try to predict outcomes from mixed fields,
which are based on available clinical response data. Most notably, an approach
based on the local effect model (LEM), where macroscopic response data in the
form of dose–effect curves is used to quantify the relative effect of the dose
delivered[20]. The parameterization, however, of the latter approach is
extensive due to the fact that every effect curve has two parameters for a
given $\alpha/\beta$–value, making the model over–parameterized. As such, it
is possible to have this model reflect the current knowledge of dose and
modality response adequately, which forms an important, albeit controversial
tool[21, 22]. Its power to predict the behavior outside of the current
knowledge therefore seems to be limited.
Cucinotta et al. attempted to incorporate the volumetric properties of the
dose deposition[23] to account for differences in track structure. They
observed that: “LET is a poor descriptor of energy deposition in small volumes
because of the diffusion of secondary electrons out of the volume and
contribution of $\delta$–rays that pass outside of the volume”. To address
this problem a quantification of the energy distribution of generated
secondary particles, or $\delta$–rays was proposed.
Such a secondary charged particle indirectly changes the behavior with respect
to the DNA damage induced. Indeed, depending on the median energy of the
spectrum the DNA damage changes accordingly if the dose is kept constant. In
the paper presented here this behavior could be easily incorporated by
considering the DNA damage for all the particles (i.e. ions and $\delta$–rays)
separately using a methodology modeled on the use of the electronic
equilibrium concept in photon cavity theory. Currently this behavior is hidden
in the $k$ parameter and it would be interesting to see if such an approach
will lead to a convergence of all $k$–values for all particles.
To take these actions fully into account an approach to provide a more
detailed model of the biological effect directly in the Monte Carlo simulation
is proposed by Sato et al.[24] This would, in theory, allow a direct
calculation of the effect in terms of energy deposited. However, as outlined
by Cucinotta this is not without problems as the behavior of low energy
electrons needs to be adequately modelled. This work predicts that the current
knowledge using the Bethe formalism, might not be suitably extended.
The results from the geometric interpretation indicate that the overall
behavior of the DNA damage induction is identical for all types of charged
particles. The only difference is in the dimensionless parameter $k$. The
latter seems to change as the ion used is heavier. Preliminary calculations
using the IMFP in water indicate that the value of $k$ diminishes as the
charged particles used are heavier (or more charged, data not shown). A
possible reason for this is that the track structure can be quite different
for different charged particles. This fact could also be an explanation for
the discrepancy found at very low energies for carbon–ions. Indeed, allowing
the parameter $k$ to be covariant with the other parameters, does provide a
more adequate fit (data nor shown).
The results for the electrons also shows a discrepancy with regard to the
generation of complex damage at lower energies. For electrons, the data on
very low energy electrons are not available in terms of energy deposition.
Indeed, the model proposed here shows a much lower incidence of complex damage
due to plasmonic effects in that region.
In summary, the model proposed here allows extension to very low energies for
electrons and protons.The fact that there are indications that the induction
of DSB’s varies linearly with dose, provides an easy implementation to dose
planning systems, given the knowledge of dose deposition spectra in a
treatment beam. An example of such implementation is provided in the Appendix.
## Appendix: Implementation in dose deposition calculations
### Mono–energetic treatment
In dose calculations a dose matrix is obtained on a dose grid Let
$\mathbf{D}=D[i,j,k]$ be the dose matrix provided. Then we can write the
amount of complex damage incurred by particles with an energy ($E$) as a
damage matrix, denoted as ($M_{cd}$). as follows:
$\displaystyle\mathbf{M_{cd}}=M_{cd}[i,j,k]$ $\displaystyle=\mathbf{D}\times
F_{cd}(E)$ (13)
$F_{cd}(E)$ then denotes a response function converting dose to damage.
### Poly–energetic treatment
Dose deposition spectra rarely consist of a field of mono–energetic electrons.
For a photon source with a given photon spectrum, an energy depositing
electron fields exists, which is roughly constant throughout the target
volume. Using Monte Carlo simulations it is possible to calculate this field
and its spectrum $\Psi(E)$. It then becomes possible to include the spectrum
in the calculation of the damage matrices. This approach has been used already
by different authors [25, 26].
$\mathbf{M_{cd}}~{}=~{}\mathbf{D}\times\frac{\int_{0}^{E_{max}}\Psi(E)F_{cd}(E)dE}{\int_{0}^{E_{max}}\Psi(E)dE}$
(14)
In the case of charged particle treatment, the particles are moderated and the
energy spectrum changes depending on the position of the point where the dose
is being deposited. It is therefore necessary to apply Equation 14 to each
point separately with knowledge of the depositing energy spectrum in that
point. Due to the closed nature of the formalism developed in this paper, it
becomes feasible to use off the shelf computing equipment.
### Application: Proton treatment
Recently, the coupling of Monte Carlo simulations in dose deposition to micro-
dosimetric code has been proposed and applied by several groups[26, 25]. Here
a two step approach is followed; 1) a general purpose Monte Carlo code (MCNPX
2.7b)[27] is used to estimate the spectrum of all different dose contributing
particles, 2) a micro dosimetric code[13] is used to determine the biological
damage.
The framework for conversion of dose to biological effect is implemented on a
simulation of a pristine 200MeV proton beam, taking into account the changing
proton spectrum. The proton simulation is performed using MCNPX. Figure 6(a)
shows the variation of the number of complex damage events as a function of
energy of the proton. In addition, the spectrum of depositing protons is shown
at a position before the Bragg peak and at the Bragg peak. In Figure 6(b) the
effect on the dose deposition is shown together with the $RBE_{cd}$ calculated
as the complex damage yield generated by the protons at that particular
position, divided by the complex damage induced by a 6MV photon beam with the
same spatial characteristics. Note, that the $RBE_{cd}$ is of the order of 1.1
with larger value of 2 a few mm distal from the Bragg peak. This is
commensurate with the cell data reported by Paganetti et al.[28] and Chaudhary
et al.[29], who showed that the radiobiological effect at the distal end of a
spread out bragg peak increases, a fact predicted by Goitein[30]. Currently,
data of direct measurement of DNA–damage in–vitro along a proton beam are
scarce. The advent of $\gamma$–H2AX measurements, as a marker for DSB–damage
is promising in this regard and has been used to investigate anti–protons[31].
(a) Impact of the spectrum (b) Proton dose deposition
Figure 6: In Fig 6(a) the spectrum at the beginning of the Bragg peak (scaled
by 0.1) is completely within the low LET regimen. While at the end of the
Bragg peak a significant part of the dose depositing protons exhibits high LET
characteristics. Fig.6(b) shows the $\mathrm{RBE_{cd}}$ (red line) together
with the damage induced by a mono–energetic proton beam.
## Acknowledgements
I am indebted to Rob Stewart, not only for generously providing the MCDS–code
for everyone to use, but also for providing much needed input on the
biological aspects of the MCDS implementation. Both Ricardo Raabe and Mike
Partridge reviewed some of the physics and made it a more rigorous paper.
Sandra Nuyts and Dirk De Ruysscher’s clinical input was also highly
appreciated.
## References
* 1. Hall EJ (1978) LET and RBE. In: Radiobiology for the radiologist, Philadelphia: Harper & Row, Publishers.
* 2. Caldecott KW (2008) Single-strand break repair and genetic disease. Nat Rev Genet 9: 619–631.
* 3. Ward JF (1985) Biochemistry of dna lesions. Radiation Research Supplement 8: pp. S103-S111.
* 4. Bethe H (1930) Zur theorie des durchgangs schneller korpuskularstrahlen durch materie. Ann Phys 5.
* 5. Barkas W, Dyer J, Heckman H (1963) Resolution of the $\sigma$–mass anomaly. Phys Rev Lett 11: 26.
* 6. Brenner DJ, Ward J (1992) Constraints on energy deposition and target size of multiply damaged sites associated with dna double-strand breaks. International Journal of Radiation Biology 61: 737-748.
* 7. ICRU, International Commission on Radiation Units and Measurements (1970) Linear energy transfer. ICRU Report 16.
* 8. ICRU, International Commission on Radiation Units and Measurements (1971) Radiation quantities and units. ICRU Report 19.
* 9. ICRU, International Commission on Radiation Units and Measurements (1983) Microdosimetry. ICRU Report 36.
* 10. Nikjoo H, O’Neill P, Goodhead DT, Terrissol M (1997) Computational modelling of low-energy electron-induced DNA damage by early physical and chemical events. Int J Radiat Biol 71: 467-473.
* 11. Stewart RD, Yu VK, Georgakilas AG, Koumenis C, Park JH, et al. (2011) Effects of radiation quality and oxygen on clustered DNA lesions and cell death. Radiation Research 176: 587-602.
* 12. Cauchy A (1853) Sur les résultats moyens d’observations de même nature, et sur les résultats les plus probables. Comptes Rendus de l’Académie des Sciences : 198-206.
* 13. Semenenko V, Stewart R (2006) Fast Monte Carlo simulation of DNA damage formed by electrons and light ions. Phys Med Biol 51: 1693-1706.
* 14. Breit G (1959) Handbuch der Physik XLI/1. Berlin, Heidelberg: Springer.
* 15. Powell C, Jablonski A (2010) Nist standard reference database 71. In: NIST Electron Inelastic–Mean–Free–Path Database — Version1.2, NIST Gaithersburg, MD.
* 16. Zhen-Yu T, Yue-Yuan X, Ming-Wen Z, Xiang-Dong L (2010) Proton inelastic mean free path in a group of organic materials in 0.05-10mev range. Chinese Physics Letters 27: 113403.
* 17. Sinden R (1994) DNA structure and function. San Diego: Academic Press.
* 18. Tanuma S, Powell CJ, Penn DR (2005) Calculations of electron inelastic mean free paths. Surface and Interface Analysis 37: 1–14.
* 19. Ziaja B, London RA, Hajdu J (2006) Ionization by impact electrons in solids: Electron mean free path fitted over a wide energy range. Journal of Applied Physics 99: 033514.
* 20. Kraemer M, Scholz M (2006) Rapid calculation of biological effects in ion radiotherapy. Phys Med Biol 51: 1959-1970.
* 21. Katz R (2003) The parameter-free track structure model of Scholz and Kraft for heavy-ion cross sections. Radiation Research 160: 724–728.
* 22. Paganetti H, Goitein M (2001) Biophysical modelling of proton radiation effects based on amorphous track models. International Journal of Radiation Biology 77: 911-928.
* 23. Cucinotta FA, Nikjoo H, Goodhead DT (2000) Model for radial dependence of frequency distributions for energy imparted in nanometer volumes from hze particles. Radiation Research 153: 459–468.
* 24. Sato T, Kase Y, Watanabe R, Niita K, Sihver L (2009) Biological dose estimation for charged-particle therapy using an improved phits code coupled with a microdosimetric kinetic model. Radiation Research 171: 107–117.
* 25. Van den Heuvel F, Locquet JP, Nuyts S (2010) Beam energy considerations for gold nano-particle enhanced radiation treatment. Physics in Medicine and Biology 55: 4509.
* 26. Hsiao Y, Stewart RD (2008) Monte carlo simulation of DNA damage induction by X-rays and selected radioisotopes. Phys Med Biol 53: 233.
* 27. Waters LS, McKinney GW, Durkee JW, Fensin ML, Hendricks JS, et al. (2007) The MCNPX Monte Carlo radiation transport code. In: Albrow, M and Raja, R, editor, Hadronic Shower Simulation Workshop. Amer Inst Physics, volume 896 of _AIP conference proceedings_ , pp. 81-90. Hadronic Shower Simulation Workshop, Batavia, IL, SEP 06-08, 2006.
* 28. Paganetti H, Niemierko A, Ancukiewicz M, Gerweck LE, Goitein M, et al. (2002) Relative biological effectiveness (RBE) values for proton beam therapy. International journal of radiation oncology, biology, physics 53: 407–421.
* 29. Chaudhary P, Marshall TI, Perozziello FM, Manti L, Currell FJ, et al. (2014) Relative biological effectiveness variation along monoenergetic and modulated bragg peaks of a 62-mev therapeutic proton beam: A preclinical assessment. Int J Radiat Oncol Biol Phys 90: 27–35.
* 30. Goitein M (2008) Radiation Oncology: A Physicist’s-Eye View. Springer Verlag.
* 31. Kavanagh JN, Currell FJ, Timson DJ, Savage KI, Richard DJ, et al. (2013) Antiproton induced dna damage: proton like in flight, carbon-ion like near rest. Sci Rep 3: –.
|
arxiv-papers
| 2013-09-29T09:53:41 |
2024-09-04T02:49:51.663462
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Frank Van den Heuvel",
"submitter": "Frank Van den Heuvel",
"url": "https://arxiv.org/abs/1309.7563"
}
|
1309.7676
|
# An upper bound on prototype set size for condensed nearest neighbor
Eric Christiansen
###### Abstract
The condensed nearest neighbor (CNN) algorithm is a heuristic for reducing the
number of prototypical points stored by a nearest neighbor classifier, while
keeping the classification rule given by the reduced prototypical set
consistent with the full set. I present an upper bound on the number of
prototypical points accumulated by CNN. The bound originates in a bound on the
number of times the decision rule is updated during training in the multiclass
perceptron algorithm, and thus is independent of training set size.
## 1 Introduction
The nearest neighbor (NN) rule assigns to an unclassified point the class of a
closest point from a set of prototypical points. The NN algorithm stores every
training point as a prototypical point and classifies new points according to
the NN rule. A nice property is that, for arbitrary class distributions, as
the number of training points goes to infinity, the error of the rule produced
by the NN algorithm converges to within twice the Bayes error [4].
Unfortunately, storing every training point as a prototypical point can be
impractical for huge training sets in terms of both memory complexity and the
time complexity of classifying according to the NN rule. As a result, many
techniques exist for reducing the size of the set of prototypical points. See
[8] and [6] for an overview. It is suggested in [1] that using a smaller set
with the NN rule may be preferable to using the entire training set, because
the VC dimensionality of a NN classifier is given by the number of its
prototypical points. However, experimental results in [8] suggest
generalization performance is best when the entire training set is used,
illustrating a tradeoff between computational complexity and generalization
performance.
This paper focuses on an algorithm for finding a consistent subset of the
training data, where a set of prototypes is consistent if it classifies the
training set correctly using the NN rule. Finding a consistent subset of
minimum cardinality is intractable [7], but several heuristic approaches
exist. One approach is the condensed nearest neighbor (CNN) algorithm, a
simple method introduced by Hart [5] that has met with empirical success
[1][8]. In this paper, I point out a striking similarity between the CNN
algorithm and the multiclass perceptron algorithm described by Collins in [3].
I develop that connection to derive an upper bound on the number of
prototypical points accumulated by CNN. The existence of this bound may help
explain CNN’s success in [1] and [8].
## 2 Condensed nearest neighbor
We will first need some definitions. Let $T=\\{({\bf x}_{i},c_{i})\\}$ be a
set of classified points, where each point ${\bf
x}_{i}\in\mbox{${\mathbb{R}}$}^{d}$ and each label $c_{i}$ comes from a finite
set of classes $C$. Let $N_{T}({\bf x})\in\mbox{${\mathbb{R}}$}^{d}$ be the
point part of the element in $T$ minimizing Euclidean distance to ${\bf x}$;
so $N_{T}({\bf x})$ is ${\bf x}$’s nearest neighbor in $T$.111To simplify
analysis, we shall assume $N_{T}$ is always well defined, meaning that every
point we consider has exactly one nearest neighbor in $T$. We will also assume
there do not exist $({\bf x},c),({\bf x}^{\prime},c^{\prime})\in T$ such that
${\bf x}={\bf x}^{\prime}$ and $c\neq c^{\prime}$; in other words, there are
no conflicting labels. Let $C_{T}({\bf x})\in C$ be the class of $N_{T}({\bf
x})$. If $T=\emptyset$, then neither $N_{T}$ nor $C_{T}$ exist, and it is
understood they cannot be said to equal anything. A set $P$ is said to be
consistent with $T$ if $C_{P}({\bf x})=c$ for every $({\bf x},c)\in T$.
Suppose $T=\\{({\bf x}_{i},c_{i})\\}$ is our training set. Then the condensed
nearest neighbor (CNN) algorithm builds a consistent subset of prototypical
points $P$ [5]. Its definition is given in Algorithm 1. Here the for loop
selects points in arbitrary order, but without repetition, from $T$.
Algorithm 1 Condensed nearest neighbor algorithm
$P\leftarrow\emptyset$
flag $\leftarrow$ True
while flag is True do
flag $\leftarrow$ False
for $({\bf x},c)$ in $T$ do
if $C_{P}({\bf x})\neq c$ then
$P\leftarrow P\cup\\{({\bf x},c)\\}$
flag $\leftarrow$ True
end if
end for
end while
For finite $T$, this algorithm is guaranteed to terminate; in the worst case,
it stops after adding every element of $T$ to $P$. It is also easy to see that
after termination, $P$ will be consistent with $T$. In practice, CNN often
finds a prototypical set far smaller than the original training set [1][8],
though this improvement is obviously more pronounced when the training set is
much larger than it needs to be.
CNN can be naturally modified for the online context; the while loop is
dropped and it is understood the for loop becomes an iteration over an
infinite stream of training data. However, in the case of overlapping class
distributions (non-zero Bayes’ error), the expected size of $P$ will grow
linearly with the number of observed points, perhaps making CNN unsuitable.
## 3 Multiclass perceptron
The multiclass perceptron (MP) algorithm [3] is listed as Algorithm 2. It can
be thought of as a generalization of the kernelized perceptron algorithm ([2],
pages 192-196) to handle multiple classes. Here,
$\phi:\mbox{${\mathbb{R}}$}^{d}\times C\rightarrow H$ is an arbitrary feature
function, where $H$ is a Hilbert space. The algorithm builds a decision rule
through iterative updates to a vector ${\bf w}$. Note a nonlinear decision
rule can be obtained through the use of a nonlinear feature function.
Algorithm 2 Multiclass perceptron algorithm
${\bf w}\leftarrow{\bf 0}$
flag $\leftarrow$ True
while flag is True do
flag $\leftarrow$ False
for $({\bf x},c)$ in $T$ do
$y\leftarrow\operatornamewithlimits{argmax}_{y\in C}{\bf w}\cdot\phi({\bf
x},y)$
if $y\neq c$ then
${\bf w}\leftarrow{\bf w}+\phi({\bf x},c)-\phi({\bf x},y)$
flag $\leftarrow$ True
end if
end for
end while
I will say a set $T=\\{({\bf x}_{i},c_{i})\\}$ is separable with margin
$\delta>0$ if there exists a unit vector ${\bf w}^{*}\in H$ such that for
every $({\bf x},c)\in T$ and for every $y\in C-\\{c\\}$ we have ${\bf
w}^{*}\cdot\phi({\bf x},c)-{\bf w}^{*}\cdot\phi({\bf x},y)\geq\delta$. We will
say $T$ has radius at most $R$ if for every $({\bf x},c)\in T$ and for every
$y\in C-\\{c\\}$ we have $||\phi({\bf x},c)-\phi({\bf x},y)||\leq R$. Note
both margin and radius depend on the feature function $\phi$ as well as the
data $T$. From [3] we have Theorem 1.
###### Theorem 1.
Let $T$ be separable with margin $\delta$ and radius at most $R$. Then the MP
algorithm updates ${\bf w}$ at most $R^{2}/\delta^{2}$ times.
Note the number of updates to ${\bf w}$ is the same as the number of points
$({\bf x},c)\in T$ that are misclassified during the construction of the
decision rule. Also note if $T$ is not separable, then it will not have a
margin $\delta>0$. In fact, the MP algorithm will never terminate.
## 4 An upper bound on prototype set size for condensed nearest neighbor
Here I develop the connection between the condensed nearest neighbor (CNN) and
multiclass perceptron (MP) algorithms. The motivating insight comes from
support vector machine (SVM) lore, which asserts that an SVM with a Gaussian
kernel is like a smoothed nearest neighbor (NN) classifier. In a similar vein,
I will show that for certain $\phi$s and ${\bf w}$s, we can use ${\bf
w}\cdot\phi({\bf x},y)$ as a proxy for the nearness of ${\bf x}$ and its
nearest neighbor of class $y$, making $\operatornamewithlimits{argmax}_{y\in
C}{\bf w}\cdot\phi({\bf x},y)$ the class of ${\bf x}^{\prime}s$ nearest
neighbor.
First we need some definitions. Let $T=\\{({\bf x}_{i},c_{i})\\}$ be a
training set, and let $P\subseteq T$. Let
$o:\mbox{${\mathbb{R}}$}^{d}\rightarrow C$ be any multifunction, where for
every $({\bf x},c)\in T$, we have $o({\bf x})\neq c$. So $o({\bf x})$ can be
any class, except the true class of ${\bf x}$. I say the set of all ${\bf w}$s
of the form given in (1) are the restricted ${\bf w}$s corresponding to $P$.
$\displaystyle{\bf w}$ $\displaystyle=\sum_{({\bf x},c)\in P}\phi({\bf
x},c)-\phi({\bf x},o({\bf x}))$ (1)
Note in (1), ${\bf w}$ has exactly the form one might expect the ${\bf w}$
from the MP algorithm to have after termination, except that every summand in
a restricted ${\bf w}$ has the unitary coefficient.
Now I define a condition on $\phi$ that makes
$\operatornamewithlimits{argmax}_{y\in C}{\bf w}\cdot\phi({\bf x},y)$ behave
like the NN rule. I will say $\phi$ is neighborly (with respect to $T$) if for
every $P\subseteq T$ and every corresponding restricted ${\bf w}$, as well as
for every training point $({\bf x},c)\in T$, we have
$\operatornamewithlimits{argmax}_{y\in C}{\bf w}\cdot\phi({\bf
x},y)=C_{P}({\bf x})$. When $|T|$ is finite, we can always find uncountably
many neighborly $\phi$s; see Appendix A for a construction using the Gaussian
kernel.
By Lemma 1, when $\phi$ is neighborly, the MP algorithm will never misclassify
the same point twice. So we will only need to consider restricted ${\bf w}$s
when analyzing the MP algorithm.
###### Lemma 1.
Let $\phi$ be neighborly. Then at all times during the execution of Algorithm
2, ${\bf w}$ will be restricted.
###### Proof.
This is a proof by induction. Note ${\bf w}$ is initially restricted. Further,
if ${\bf w}$ is restricted before the $n^{\text{th}}$ update, then it will be
restricted after the $n^{\text{th}}$ update. For let $({\bf x},c)\in T$ be the
point causing the $n^{\text{th}}$ update. So
$\operatornamewithlimits{argmax}_{y\in C}{\bf w}\cdot\phi({\bf x},y)\neq c$.
But $\phi$ is neighborly, and ${\bf w}$ is restricted before this update, so
$\operatornamewithlimits{argmax}_{y\in C}{\bf w}\cdot\phi({\bf
x},y)=C_{P}({\bf x})$. So $C_{P}({\bf x})\neq c$. So $({\bf x},c)\not\in P$,
and thus ${\bf w}$ is restricted after the $n^{\text{th}}$ update. ∎
The CNN algorithm is essentially the MP algorithm 2 in the special case $\phi$
is neighborly, so we should be able to apply to the CNN algorithm theorems
that apply to the MP algorithm. In that spirit, Lemma 2 says we can apply
Theorem 1 to the CNN algorithm.
###### Lemma 2.
Let $\phi$ be neighborly, and suppose $T$ is separable with margin $\delta$
and radius at most $R$. Then the CNN algorithm updates $P$ at most
$R^{2}/\delta^{2}$ times.
###### Proof.
The idea is to morph the MP algorithm into the CNN algorithm without changing
the number of misclassifications, which are in one-to-one correspondence with
the number of updates to ${\bf w}$ in the MP algorithm, and are in one-to-one
correspondence with the number of updates to $P$ in the CNN algorithm.
We begin by inserting two lines of code into the MP algorithm, obtaining
Algorithm 3. The new lines in Algorithm 3 are lines 2 and 10. These lines
clearly have no effect on the number of updates to ${\bf w}$, but note in
Algorithm 3 the number of updates to $P$ equals the number of updates to ${\bf
w}$.
Algorithm 3 Hybrid multiclass perceptron algorithm
1: ${\bf w}\leftarrow{\bf 0}$
2: $P\leftarrow\emptyset$
3: flag $\leftarrow$ True
4: while flag is True do
5: flag $\leftarrow$ False
6: for $({\bf x}_{i},c_{i})$ in $T$ do
7: $y_{i}\leftarrow\operatornamewithlimits{argmax}_{y\in C}{\bf
w}\cdot\phi({\bf x}_{i},y)$
8: if $y_{i}\neq c_{i}$ then
9: ${\bf w}\leftarrow{\bf w}+\phi({\bf x}_{i},c_{i})-\phi({\bf x}_{i},y_{i})$
10: $P\leftarrow P\cup\\{({\bf x}_{i},c_{i})\\}$
11: flag $\leftarrow$ True
12: end if
13: end for
14: end while
We then use Lemma 1, together with the assumption $\phi$ is neighborly, to
replace $\operatornamewithlimits{argmax}_{y\in C}{\bf w}\cdot\phi({\bf
x}_{i},y)$ with $C_{P}({\bf x}_{i})$ in Algorithm 3, yielding Algorithm 4.
Algorithm 4 Hybrid condensed nearest neighbor algorithm
1: ${\bf w}\leftarrow{\bf 0}$
2: $P\leftarrow\emptyset$
3: flag $\leftarrow$ True
4: while flag is True do
5: flag $\leftarrow$ False
6: for $({\bf x}_{i},c_{i})$ in $T$ do
7: if $C_{P}({\bf x}_{i})\neq c_{i}$ then
8: ${\bf w}\leftarrow{\bf w}+\phi({\bf x}_{i},c_{i})-\phi({\bf
x}_{i},C_{P}({\bf x}_{i}))$
9: $P\leftarrow P\cup\\{({\bf x}_{i},c_{i})\\}$
10: flag $\leftarrow$ True
11: end if
12: end for
13: end while
Note Algorithm 4 updates $P$ exactly as many times as the MP algorithm updated
${\bf w}$, and from Algorithm 4 we can simply delete both lines referencing
${\bf w}$ (lines 1 and 8). This yields the CNN algorithm, without changing the
number of updates to $P$. So by Theorem 1, we are done. ∎
Since Lemma 2 gives a bound for any neighborly feature function, the set of
all neighborly feature functions yields a set of bounds. It is natural to
choose the best bound, which is the essence of Theorem 2. Note the number of
updates to $P$ is exactly the ultimate size of $P$, or equivalently the number
of points accumulated by the CNN algorithm.
###### Theorem 2.
Let $\Phi=\\{\phi:\phi\text{ is neighborly}\\}$ and suppose, for each
$\phi\in\Phi$, $T$ is separable with margin $\delta_{\phi}$ and radius at most
$R_{\phi}$. Then the CNN algorithm accumulates at most
$\inf_{\phi\in\Phi}R_{\phi}^{2}/\delta_{\phi}^{2}$ prototypical points.
###### Proof.
By Lemma 2, for any neighborly $\phi$, the CNN algorithm accumulates at most
$R_{\phi}^{2}/\delta_{\phi}^{2}$ representative points. Thus since $\Phi$ is
the set of all neighborly $\phi$s, the CNN algorithm accumulates at most
$\inf_{\phi\in\Phi}R_{\phi}^{2}/\delta_{\phi}^{2}$ representative points. ∎
So $|P|$ is bounded above by the best bound over all neighborly $\phi$s.
## 5 Conclusion
In Theorem 2 I presented a bound on the number of prototypical points
accumulated by the condensed nearest neighbor (CNN) algorithm. This bound came
from a bound on the number of updates to the decision rule in the multiclass
perceptron algorithm. Unfortunately, as with the multiclass perceptron bound
presented in [3], estimating the bound is likely to be too expensive to be
practical. Fortunately, as with the multiclass perceptron bound, this bound is
independent of the size of the training set used as input to the CNN
algorithm. Perhaps the existence of this bound may help explain CNN’s
empirical success.
## Appendix A Neighborly $\phi$s
Let $T=\\{({\bf x}_{i},c_{i})\\}$ be a finite training set, where ${\bf
x}_{i}\in\mbox{${\mathbb{R}}$}^{d}$, $c_{i}\in C$, and $|C|$ is finite. I will
show there exist uncountably many neighborly feature functions
$\phi:\mbox{${\mathbb{R}}$}^{d}\times C\rightarrow H$, where $H$ is a Hilbert
space. In other words, there exist uncountably many functions
$\phi:\mbox{${\mathbb{R}}$}^{d}\times C\rightarrow H$ such that for any
$P\subseteq T$ and restricted ${\bf w}$ corresponding to $P$, as well as for
any $({\bf x},c)\in T$, we have $\operatornamewithlimits{argmax}_{y\in C}{\bf
w}\cdot\phi({\bf x},y)=C_{P}({\bf x})$. The idea is to construct an
uncountable family of such functions using the Gaussian kernel.
Let $\psi_{\sigma}:\mbox{${\mathbb{R}}$}^{d}\rightarrow H$ correspond to the
Gaussian kernel, so that
$\psi_{\sigma}({\bf x})\cdot\psi_{\sigma}({\bf
x}^{\prime})=e^{\frac{-1}{2\sigma^{2}}||{\bf x}-{\bf x}^{\prime}||^{2}}.$
We know such a $\psi_{\sigma}$ exists ([2], pages 294-299). Let
$\phi_{\sigma}({\bf x},c)\in H$ be a vector such that the $i^{\text{th}}$
element $\phi_{\sigma}({\bf x},c)_{i}$ of $\phi_{\sigma}({\bf x},c)$ is
$\phi_{\sigma}({\bf x},c)_{i}=\left\\{\begin{tabular}[]{ c
l}$\psi_{\sigma}({\bf x})_{j}$&if there exists a $j$ such that
$i=(j-1)|C|+c$\\\ $0$&otherwise\\\ \end{tabular}\right..$
The idea behind the above definition is to make
$\phi_{\sigma}({\bf x},c)\cdot\phi_{\sigma}({\bf
x}^{\prime},c^{\prime})=I(c=c^{\prime})\psi_{\sigma}({\bf
x})\cdot\psi_{\sigma}({\bf x}^{\prime}),$
where $I$ is the indicator function on Boolean inputs given by
$I(\texttt{true})=1$ and $I({\texttt{false}})=0$.
I claim $\phi_{\sigma}$ will be neighborly if we choose $\sigma>0$
sufficiently small. First we need the following result.
###### Lemma 3.
Let $P\subseteq T$ and let ${\bf w}$ be restricted corresponding to $P$. Let
$({\bf x}^{\prime},c^{\prime})\in T$. Then there exists a $\sigma^{*}>0$ such
that for every $0<\sigma<\sigma^{*}$, we have
$\operatornamewithlimits{argmax}_{y\in C}{\bf w}\cdot\phi_{\sigma}({\bf
x}^{\prime},y)=C_{P}({\bf x}^{\prime})$.
###### Proof.
Note ${\bf w}$ is of the form
$\displaystyle{\bf w}$ $\displaystyle=\sum_{({\bf x},c)\in
P}\phi_{\sigma}({\bf x},c)-\phi_{\sigma}({\bf x},o({\bf x})).$
So
$\displaystyle{\bf w}\cdot\phi_{\sigma}({\bf x}^{\prime},y)$
$\displaystyle=\sum_{({\bf x},c)\in P}I(c=y)\psi_{\sigma}({\bf
x})\cdot\psi_{\sigma}({\bf x}^{\prime})-I(o({\bf x})=y)\psi_{\sigma}({\bf
x})\cdot\psi_{\sigma}({\bf x}^{\prime})$ (2) $\displaystyle=\sum_{({\bf
x},c)\in P}\left[I(c=y)-I(o({\bf x})=y)\right]e^{\frac{-1}{2\sigma^{2}}||{\bf
x}-{\bf x^{\prime}}||^{2}}.$ (3)
I want to show that as $\sigma\rightarrow 0$, the summation in (3) will be
dominated by the term corresponding to the nearest neighbor of ${\bf
x}^{\prime}$ in $P$. The intuition behind this is that the summands of (3)
correspond to individual Gaussians in a mixture of Gaussians, where the
typical constraints on the mixing coefficients are dropped. For large
$\sigma$, the value of the mixture at a point might depend significantly on
several nearby Gaussians, but as we uniformly shrink the variance of the
Gaussians, the value of the mixture at a point will depend almost entirely on
the Gaussian nearest that point, as is borne out in the following analysis.
First, I divide through by $e^{\frac{-1}{2\sigma^{2}}||N_{P}({\bf
x}^{\prime})-{\bf x^{\prime}}||^{2}}$, giving us
$\displaystyle(\ref{eqn:appndxdom})$ $\displaystyle\propto\sum_{({\bf x},c)\in
P}\left[I(c=y)-I(o({\bf
x})=y)\right]e^{\frac{1}{2\sigma^{2}}\left[||N_{P}({\bf x}^{\prime})-{\bf
x^{\prime}}||^{2}-||{\bf x}-{\bf x^{\prime}}||^{2}\right]}$ (4)
$\displaystyle=\left[I(C_{P}({\bf x}^{\prime})=y)-I(o(N_{P}({\bf
x}^{\prime}))=y)\right]$ $\displaystyle\quad+\sum_{({\bf x},c)\in P:{\bf
x}\neq N_{P}({\bf x}^{\prime})}\left[I(c=y)-I(o({\bf
x})=y)\right]e^{\frac{1}{2\sigma^{2}}\left[||N_{P}({\bf x}^{\prime})-{\bf
x^{\prime}}||^{2}-||{\bf x}-{\bf x^{\prime}}||^{2}\right]}.$ (5)
Let $S$ be the summation given in (5). Note for every $({\bf x},c)\in P$ such
that ${\bf x}\neq N_{P}({\bf x}^{\prime})$, we have
$||N_{P}({\bf x}^{\prime})-{\bf x^{\prime}}||^{2}-||{\bf x}-{\bf
x^{\prime}}||^{2}<0.$
Thus as $\sigma\rightarrow 0$, every term in S will go to zero. So by the
definition of a limit, there is a $\sigma^{*}>0$ such that $\sigma<\sigma^{*}$
implies $|S|<\frac{1}{2}$. Let $\sigma<\sigma^{*}$. Then if $y=C_{P}({\bf
x}^{\prime})$, we have $(\ref{eqn:appndxsep})\geq 1-|S|>\frac{1}{2}$. If
$y\neq C_{P}({\bf x}^{\prime})$, we have $(\ref{eqn:appndxsep})\leq
0+|S|<\frac{1}{2}$. So $\operatornamewithlimits{argmax}_{y\in C}{\bf
w}\cdot\phi_{\sigma}({\bf x}^{\prime},y)=C_{P}({\bf x}^{\prime})$. ∎
With Lemma 3 I can prove Theorem 3, establishing the desired result.
###### Theorem 3.
For any finite $T$, there exist uncountably many neighborly feature functions.
###### Proof.
By Lemma 3, for every $P\subseteq T$, restricted ${\bf w}$, and $({\bf
x}^{\prime},c)\in T$, we can find a $\sigma^{*}>0$ such that
$\sigma<\sigma^{*}$ implies $\operatornamewithlimits{argmax}_{y\in C}{\bf
w}\cdot\phi_{\sigma}({\bf x}^{\prime},y)=C_{P}({\bf x}^{\prime})$. Let
$\Sigma^{*}$ be the set of all such $\sigma^{*}$s. Since there are finitely
many subsets $P$ of $T$, finitely many restricted ${\bf w}$s corresponding to
any given $P$, and finitely many elements in $T$, the set $\Sigma^{*}$ is
finite and thus has a minimum element $\sigma^{\text{min}}$.
Thus for any of the uncountably many $\sigma$s such that
$\sigma<\sigma^{\text{min}}$, there exists a neighborly feature function
$\phi_{\sigma}$. Thus the set of all neighborly feature functions contains a
subset of uncountable cardinality, and is thus uncountable. ∎
## Acknowledgements
I would like to thank Mehran Bozorgi, Charles Elkan, and Matus Telgarsky for
acting as sounding boards for this idea, and I would like to thank Charles
Elkan and Nicolaus Hepler for constructively criticizing the drafts of this
paper.
This work was funded by NSF grant #SBE-0542013 to the Temporal Dynamics of
Learning Center, an NSF Science of Learning Center.
## References
* [1] Fabrizio Angiulli. Fast condensed nearest neighbor rule. In ICML ’05: Proceedings of the 22nd international conference on Machine learning, pages 25–32, New York, NY, USA, 2005. ACM.
* [2] Christopher M. Bishop. Pattern Recognition and Machine Learning (Information Science and Statistics). Springer, August 2006.
* [3] Michael Collins. Discriminative training methods for hidden Markov models: theory and experiments with perceptron algorithms. In EMNLP ’02: Proceedings of the ACL-02 conference on empirical methods in natural language processing, pages 1–8, Morristown, NJ, USA, 2002\. Association for Computational Linguistics.
* [4] Luc Devroye. On the inequality of Cover and Hart in nearest neighbor discrimination. IEEE Transactions on Pattern Analysis and Machine Intelligence, 3:75–78, 1981.
* [5] P. Hart. The condensed nearest neighbor rule (corresp.). Information Theory, IEEE Transactions on, 14(3):515–516, May 1968\.
* [6] Godfried Toussaint. Proximity graphs for nearest neighbor decision rules: recent progress. In Proceedings of the 34th Symposium on the INTERFACE, pages 17–20, 2002.
* [7] Gordon Wilfong. Nearest neighbor problems. In SCG ’91: Proceedings of the seventh annual symposium on Computational geometry, pages 224–233, New York, NY, USA, 1991. ACM.
* [8] D. Randall Wilson and Tony R. Martinez. Reduction techniques for instance-based learning algorithms. In Machine Learning, pages 257–286, 2000.
|
arxiv-papers
| 2013-09-29T23:45:59 |
2024-09-04T02:49:51.682533
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Eric Christiansen",
"submitter": "Eric Christiansen",
"url": "https://arxiv.org/abs/1309.7676"
}
|
1309.7740
|
# The quantum approach to human reasoning does explain the belief-bias effect
E. D. Vol [email protected] B. Verkin Institute for Low Temperature Physics
and Engineering of the National Academy of Sciences of Ukraine 47, Lenin Ave.,
Kharkov 61103, Ukraine.
###### Abstract
Based on the ideas of quantum physics and dual-process theory of human
reasoning that takes into account two primary mechanisms of reasoning : 1)
deductive rational thinking and 2) intuitive heuristic judgment, we proposed
the ”quantum” approach to practical human logic that allows one to specify the
most distinctive peculiarities in activity of two reasoning systems mentioned
above and in addition to describe phenomenologically well-established
experimentally belief-bias effect .
###### pacs:
05.40.-a
## I Introduction
The idea that some essential human values and concepts may be incompatable
with each other had originated long before the beginning of scientific
psychology. By distinct ways this idea was justified by such outstanding
philosophers and thinkers as G.W. Leibniz, N.Machiavelli and I. Kant. The
interested reader can find detail account of the history of this idea with
relevant references in 1s . However only in the XX century with the rise of
quantum theory this idea has received adequate scientific expression in the
language of the Bohr’s Complementarity Principle. We shall give here only two
distinctive quotations of founding fathers of quantum mechanics that are
clearly demonstrating their profound understanding of the inconsistancy of
some basic concepts relating to human psychology. So, in the paper of 1948 ”On
the Notions of Causality and Complementarity” N. Bohr wrote:”Recognition of
complementary relationship is not least required in psychology, where the
conditions for analysis and synthesis of experience exhibit striking analogy
with the situation in atomic physis. In fact, the use of words like thoughts,
and sentiments, equally indispensable to illustrate the diversity of psychical
experience, pertain to mutually exclusive situations characterized by a
different drawing of the line of separation between subject and object. In
particular, the place left for the feeling of volition is afforded by the very
circumstance that situations where we experience freedom of will are
incompatible with psychological situations where causal analysis is reasonably
attempted. In other words, when we use the phrase ”I will” we renounce
explanatory argumentation. In fact, the use which we make of words like
”thought” and ”feeling,” or ”instinct” and ”reason” to describe psychic
experiences of different types, shows the existence of characteristic
relationships of complementarity conditioned by the peculiarity of
introspection” 2s . On the other hand W.Pauli drew particular attention to the
problem of relation between complementarity of mental and physical aspects of
the same reality. In his inspiring paper ”The influence of archetypal ideas on
the scientific theories of Kepler” 2s he wrote:” “The general problem of the
relationship between psyche and physics, between inside and outside, will
hardly be solved with the notion of a ‘psychophysical parallelism, put forward
in the past century. However, modern science has perhaps brought us closer to
a more satisfying conception of this relationship insofar as it introduced the
concept of complementarity within physics. It would be most satisfactory if
physis and psyche could be conceived as complementary aspects of the same
reality.”.Unfortunately at the time these deep ideas are not influenced the
development of experimental psychology. All the more remarkable that modern
cognitive psychology irrespectively came in fact to the similar conclusions.In
particular numerous experts in so different areas of cognitive psychology as
attention,memory,decision making, learning with one accord believe that dual
processes and dual systems play fundamental role for nearly all basic
cognitive mechanisms in human mind (see e.g. 4s for brief review of dual-
process theory in reasoning with the list of necessary references).In what
follows we are interested only in human reasoning where two primary dual
systems of interest can be specified. One of these systems we will call it
below as deductive reasoning system (DRS) is rational, sequental and
consistent but acts relatively slow while the other - we will call it further
as heuristic reasoning system (HRS) is intuitive, fast, automatic, but to a
large extent influenced by emotions and last unconsious experience.Numerous
researches and experiments conclusively proved that there is hidden
interaction between these two cognitive systems such that a reasoning subject
is not aware of this.The belief-bias effect is the most striking manifestation
of such interaction. Roughly speaking the belief-bias effect is the innate
tendency of reasoning subjects to be more likely to accept reasons ang
arguments if they find them believable and to pay less attention of their
logical validity. The main goal of present paper is based on quantum ideas of
complementarity and dual-process theory in human reasoning to describe the
belief-bias effect phenomenologically by purely logical tools. To this end we
will use also the simplified version of discrete-continuous logic that was
formulated earlier in author preprint 5s .
The remainder of the paper is organized as follows. In chapter 2 we briefly
remind basic facts relating to discrete-continuous logic (DCL) that are
necessary for the understanding of the present paper.The main contribution of
this chapter is the interpretation of the general propositions in DCL as the
integral mental structures that consist both of logical and heuristic
constituents. Under such interpretation these two constituents of the
proposition can be considered as complementary to each other exactly like two
noncommuting observables in quantum mechanics.In chapter 3 we state the
uncertainty relation that just reflects the complementary nature of such
concepts as logical rigour and the heuristic grasp. And finally in chapter 4
using only logical tools we give the phenomenological explanation of the
belief-bias effect. Now let us go to the details.
## II Preliminiries
In this part we briefly remind for the reader convenience the necessary facts
relating to the discrete -continuous logic that were outlined more detail in
author preprint 5s . So, we will consider as the primary objects of our study
the set of general propositions (GP) -$\left\\{A_{j}\right\\}_{\text{ }}$that
may be represented by $2\times 2$ positive definite matrices with unit trace
of the following form:
$A_{j}=\begin{pmatrix}p_{j}&i\alpha_{j}\\\ -i\alpha_{j}&1-p_{j}\end{pmatrix},$
(1)
(where $i$ is imaginary unit). In this case the negation of such proposition
\- ($not$ $A_{j}$) may be defined as $(not$
$A)=\begin{pmatrix}1-p_{j}&-i\alpha_{j}\\\ i\alpha_{j}&p_{j}\end{pmatrix}$. It
turns out that in addition to negation another but already two place operation
-$\bigtriangleup$ (which is the analogue of strong disjunction in ordinary
Boolean logic) can be introduced in DCL according to the next definition:
Let $A=\begin{pmatrix}p&i\alpha\\\ -i\alpha&1-p\end{pmatrix}$ and
$B=\begin{pmatrix}q&i\beta\\\ -i\beta&1-q\end{pmatrix}$ then
$(A\bigtriangleup B)=\begin{pmatrix}R&i\gamma\\\ -i\gamma&1-R\end{pmatrix},$
(2)
where $R=p+q-2pq+2\alpha\beta$,
$\gamma=\alpha\left(1-2q\right)+\beta\left(1-2p\right)$. Comparing
representation Eq. (1) with standard form of density matrix of the mixed state
of two-level quantum system that looks as
$\rho=\begin{pmatrix}\frac{1+P_{z}}{2}&\frac{P_{x}-iP_{y}}{2}\\\
\frac{P_{x}+iP_{y}}{2}&\frac{1-P_{z}}{2}\end{pmatrix}$ (where
$P=(P_{x},P_{y},P_{z})$ is the Bloch vector of the state ) we see that GP may
be represented by the similar way but in this case $x-$ component of the Bloch
vector is equal to zero. In the rest of the paper we will use such reduced
Bloch representation for the arbitrary proposition $A,$ that is: $\
A=\begin{pmatrix}\frac{1+P_{z}}{2}&\frac{-iP_{y}}{2}\\\
\frac{iP_{y}}{2}&\frac{1-P_{z}}{2}\end{pmatrix}$ with $P=(P_{y},P_{z})$. In
this case it is convenient to introduce the complex vector $P=P_{z}-iP_{y}$
which we call further as representating vector (RV) of proposition $A$. It is
easy to verify directly that the RV of proposition $(notA)$ is equal to $(-P)$
and RV of proposition $(A\bigtriangleup B)$ is equal to $(-PQ)$ (where $Q$ is
RV of $B$ ).Note also the useful relation connecting negation with operation
$\bigtriangleup$:
$not(A\bigtriangleup
B)\equiv(A\overline{\bigtriangleup}B)=(notA)\bigtriangleup
B=A\bigtriangleup(notB).$ (3)
It should be noted that unlike of ordinary Boolean logic in DCL it is possible
to define the whole one-parameter group of continuous logical operations
(logical rotations of propositions in the plane $P_{y}-P_{z}$) according to
the following rule: if proposition $A$ has the RV -$P$ then rotated at an
angle $\Phi$ proposition $A^{{{}^{1}}}$ has RV - $P{{}^{1}}$ with components:
$\displaystyle P_{y}^{{{}^{1}}}$ $\displaystyle=$ $\displaystyle
P_{y}\cos\phi+P_{z}\sin\Phi$ $\displaystyle P_{z}^{{{}^{1}}}$ $\displaystyle=$
$\displaystyle P_{z}\cos\Phi-P_{y}\sin\Phi$ (4)
It is easy to see that the negation of any proposition coincides with logical
rotation of it at an angle $\pi$ and in addition that if one rotates the GP
$A$ at an angle $\Phi_{1}$ and the other proposition $B$ at an angle
$\Phi_{2}$ then the proposition $(A\bigtriangleup B)$ will be rotated at an
angle $\Phi_{1}+\Phi_{2}$. Thus all logical operations in DCL obtain quite
clear geometric meaning. Now after describing the syntax of DCL we can pass to
the more difficult task: clarification of its semantics that is the
interpretation both the meaning of general propositions and logical operations
with them. It should be noted that interpretation that we are going to propose
here is not the only possible but it is appropriate for our ultimate goal
namely to explain the beliefe-bias effect in human reasoning from pure logical
point of view. So, as before we will assume that diagonal elements of
representing matrix for arbitrary GP describes degree of its logical validity
(from DRS point of view) while its nondiagonal elements we will interpret as
the believability of the same proposition inspired by the heuristic reasoning
system (HRS). This interpretation can be expressed more precisely as follows.
Let us introduce two projection operators:$P_{1}$ and $P_{2}$
$\left(P_{1}^{2}=P_{1}\text{, }P_{2}^{2}=P_{2}\right)$ according to the
definition: $P_{1}=\frac{1+\sigma_{z}}{2}$ and $P_{2}=\frac{1+\sigma_{y}}{2}$.
It is easy to see that average values of these operators in the state whose
density matrix coincides with representating matrix of proposition
$A=\begin{pmatrix}p&i\alpha\\\ -i\alpha&1-p\end{pmatrix}$ give us the
propabilities of its logical plausability $p_{t}$ and its believability
$p_{b\text{ }}$respectively.Thus we obtain
$\displaystyle p_{t}$ $\displaystyle=$ $\displaystyle\left\langle
P_{1}\right\rangle=Sp(P_{1}A)=p$ (5) and $\displaystyle p_{c}$
$\displaystyle=$ $\displaystyle\left\langle
P_{2}\right\rangle=Sp(P_{2}A)=\frac{1-2\alpha}{2}$
In connection with above interpretation we want to point out two important
marginal GP: 1) $T=\begin{pmatrix}1&0\\\ 0&0\end{pmatrix}$-true proposition,
and 2) $B=\begin{pmatrix}\frac{1}{2}&-\frac{i}{2}\\\
\frac{i}{2}&\frac{1}{2}\end{pmatrix}$ \- highest possible believable
proposition and their negations: $F=(notT)-$false proposition and
$U=(notB)$-unbelievable proposition. Note in addition that noncommutativity of
operators $P_{1}$and $P_{2}$ implies that main predicates of arbitrary GP
(plausibility and belief) may be considered as complementary (in the sence of
quantum theory) aspects of the same proposition. This important fact implies
specific uncertainty relation for the observables $P_{1}$and $P_{2}$ connected
with any GP. The simple derivation of these relation is the subject of the
next section of the presenr paper.
## III The Uncertainty Relation between predicates plausibility and
believability in DCL.
To derive the required uncertainty relation it is convenient to represent any
GP $A$ in the Bloch form:
$A=\begin{pmatrix}\frac{1+P_{z}}{2}&\frac{-iP_{y}}{2}\\\
\frac{iP_{y}}{2}&\frac{1-P_{z}}{2}\end{pmatrix}.$According to definition the
uncertainty of logical truth for the proposition $A$ can be written with the
help of operator $P_{1}=\frac{1+\sigma_{z}}{2}$ as: $\bigtriangleup
p_{t}^{2}\equiv\overline{\left(\frac{1+\sigma_{z}}{2}-\overline{\frac{1+\sigma_{z}}{2}}\right)^{2}}=\frac{1}{4}\left(1-\overline{\sigma_{z}}^{2}\right)=\frac{1-P_{z}^{2}}{4}$.In
the similar manner the uncertainty of believability of the same proposition is
equal to : $\bigtriangleup
p_{c}=\frac{1}{4}\left(1-\overline{\sigma_{y}}^{2}\right)=\frac{1}{4}\left(1-P_{y}^{2}\right).$By
adding these two expressions we obtain: $\bigtriangleup
p_{A}^{2}\equiv\bigtriangleup p_{t}^{2}+\bigtriangleup
p_{c}^{2}=\frac{1}{4}\left(2-P_{y}^{2}-P_{z}^{2}\right).$Finally taking into
accout that $P_{y}^{2}+P_{z}^{2}\leqslant 1$ we get the desired relations:
$\frac{1}{4}\leqslant\bigtriangleup p_{A}^{2}\leqslant\frac{1}{2}.$ (6)
The notable fact should be mentioned here:if one takes two propositions $A$
and $B$ with RV $P$ and $Q$ respectively then according above calculation one
can write two equations 1) $\bigtriangleup
p_{A}^{2}=\frac{\left(2-P^{2}\right)}{4}$ and 2)$\bigtriangleup
p_{B}^{2}=\frac{\left(2-Q^{2}\right)}{4}.$
On the other hand as we marked earlier the proposition $\left(A\bigtriangleup
B\right)$ has RV $(-PQ)$ and hence its uncertainty is equal to
$\bigtriangleup$ $p_{(A\bigtriangleup
B)}^{2}=\frac{\left(2-P^{2}Q^{2}\right)}{4}$ .As long as $P^{2}Q^{2}\leqq
P^{2},Q^{2}$ one can conclude that $\bigtriangleup p_{(A\bigtriangleup
B)}^{2}\geqslant(\bigtriangleup p_{A}^{2},\bigtriangleup p_{B}^{2})$ and hence
as a final result of logical operation $\bigtriangleup$ the ending uncertainty
of proposition can only increases.We would like to hope that properly
organized experiments with specially selected reasoning tasks will be able to
confirm (or may be disprove) the proposed uncertainty relations
$\left(6\right)$.Now we come back to the main goal of present paper: the
explanation of the belief-bias effect in human reasoning.
## IV Many valued probabilistic logic and the Belief-Bias effect.
In this part we will try to describe (phenomenologically) the belief-bias
effect in human reasoning by purely logical tools.For this purpose it is
convenient to use some version of probabilistic many-valued logic that in some
sense can be considered as simplified version of original DCL.Really if in
original version of DCL we restrict ourselves only by discrete set of logical
rotations with angles: $0,\frac{2\pi}{N}...\frac{2\pi}{N}\left(N-1\right)$ we
obtain the closed logic with N marginal propositions which possess
representing matrices: $A_{0},A_{1}....A_{N-1}$ (where
$A_{0}=\begin{pmatrix}1&0\\\ 0&0\end{pmatrix},...$
$A_{k}=\begin{pmatrix}\frac{1+\cos\frac{2\pi}{N}}{2}&-\frac{i}{2}\sin\frac{2\pi}{N}\\\
\frac{i}{2}\sin\frac{2\pi}{N}&\frac{1-\cos\frac{2\pi}{N}}{2}\end{pmatrix}(k=1..N-1)$).
In the case when we are not interested in the ”quantum correlations” between
these marginal propositions we can consider them as approximately independent
quantities and associate with these propositions the logical basis consisting
of N distinct logical alternatives.Acting in this manner one can pass from
original DCL to standard many- valued probabilistic Boolean logic. After this
remark we will examine further four valued probabilistic logic every
proposition of which can be represented as $4\times 4$ diagonal matrix :
$A=diag\left(P_{1},P_{2},P_{3},P_{4}\right)$.Here we mean that the space of
these propositions is a tensor product of two spaces with $2\times 2$ diagonal
matrices, that is:
$A=\sum\limits_{i}a_{i}T_{i}\otimes B_{i},$ (7)
where $T_{i}=\begin{pmatrix}p_{i}&0\\\ 0&1-p_{i}\end{pmatrix},$
$B_{i}=\begin{pmatrix}q_{i}&0\\\ 0&1-q_{i}\end{pmatrix}$and
$\sum\limits_{i}a_{i}=1$. In addition we assume that matrices $T_{i}$ in the
decomposition Eq. (7) are associated with the activity of deductive cognitive
subsystem (DRS),while matrices $B_{i}$ are connected with its heuristic
subsystem (HRS).Thus the basis of this logic consists of four
propositions:$1)$ truth-believable $TB=diag(1,0,0,0),2)$ truth-unbelievable
$TU=diag(0,1,0,0),3)$false-believable $FB=diag(0,0,1,0)$ and $4)$false-
unbelievable $FU=diag(0,0,0,1).$Our next step is to define basic logical
operations that can be implemented with such propositions. The interpetation
that we have adopted above implies that the negation of proposition $A$ must
be defined as $\left(notA\right)=diag\left(P_{4},P_{3},P_{2},P_{1}\right)$.
The certain dilemma arises however when we want to define the conjunction of
two propositions $A=diag\left(P_{1},P_{2},P_{3},P_{4}\right)$ and
$B=diag\left(Q_{1},Q_{2},Q_{3},Q_{4}\right)$. We have proposed here the
following definition:
$C\equiv\left(AandB\right)=diag\left(C_{1},C_{2},C_{3},C_{4}\right),$ (8)
where $C_{1}=P_{1}\left(Q_{1}+Q_{2}\right)+P_{2}Q_{1}$, $C_{2}=P_{2}Q_{2}$,
$C_{3}=P_{1}\left(Q_{3}+Q_{4}\right)+P_{2}Q_{3}+P_{3}+P_{4}\left(Q_{1}+Q_{3}\right)$,
$C_{4}=P_{2}Q_{4}+P_{4}\left(Q_{2}+Q_{4}\right)$. This definition of
conjuction namely Eq. (8) certainly needs to be explained.First of all we note
that definition Eq. (8) satisfies to the necessary symmetry condition :
$\left(AandB\right)=\left(BandA\right)$as it should be.In addition if one
takes the projection of conjunction Eq. (8) in DRS (first reasoning subsystem)
the result is:$\left(AandB\right)_{1}=\begin{pmatrix}pq&0\\\
0&1-pq\end{pmatrix}\equiv\left(A_{1}andB_{1}\right)$ where $p=P_{1}+P_{2}$ and
$q=Q_{1}+Q_{2}$. This result obviously consistent with definition of
conjunction in ordinary probabilistic Boolean logic. On the other hand if one
takes the projection of Eq. (8) in HRS (second reasoning subsystem) the
obtained result reads as:
$\left(AandB\right)_{2}=\begin{pmatrix}1-\left(P_{2}+P_{4}\right)\left(Q_{2}+Q_{4}\right)&\\\
\left(P_{2}+P_{4}\right)\left(Q_{2}+Q_{4}\right)&\end{pmatrix}.$ (9)
We see that conjunction in heuristic system differs from standatd logical
conjunction .In our opinion this distinction explicitly reflects (from
phenomenological point of view) the essential difference existing between two
reasoning systems when they operate jointly. In particular the definition Eq.
(9) implies for two basic marginal propositions in second reasoning subsystem:
$B=\begin{pmatrix}1&0\\\ 0&0\end{pmatrix}$ responding to the statement of
unconditional belief and $D=\begin{pmatrix}0&0\\\ 0&1\end{pmatrix}$ -which is
the most doubtful statement, the next conjunction
relations:$\left(BandB\right)=B$, $\left(BandD\right)=\left(DandB\right)=B$
and $\left(DandD\right)=D$.Thus we obtain that the unconditional belief when
it conflicts with certain doubtful one always overcomes it. Now if one takes
the expression Eq. (8) for granted then he (she)can define another logical
operations (in particular implication that we especially interested in )
without any obstacles.To this end one should be guided by two relations of
ordinary logic which as we assume continue to be valid in our case as well: 1)
$\left(AorB\right)=not\left[\left(notA\right)and\left(notB\right)\right]$ and
2) $(A\Longrightarrow B)=\left(notA\right)orB$ .Acting in this manner we
obtain for the implication $(A\Longrightarrow B)$ the required relation:
$I\equiv(A\Longrightarrow B)=diag\left(I_{1},I_{2},I_{3},I_{4}\right),$ (10)
where $I_{1}=p_{4}\left(q_{1}+q_{3}\right)+p_{2}q_{1}$,
$I_{2}=p_{3}+q_{2}\left(1-p_{3}\right)+p_{1}q_{1}+p_{4}q_{4}$,
$I_{3}=p_{2}q_{3}$, $I_{4}=p_{1}\left(q_{3}+q_{4}\right)+p_{2}q_{4}$. The
expression Eq. (10) for the implication of two probabilistic propositions in
four- valued logic is the foundation for our following explanation of bias-
belief effect.Note that here we are going to demonstrate only the simplest
case of the application of the approach proposed. The detail quantitative
analysis of numerous possible situations connected with the interaction
between DRS and HRS will be realized by us at length in separate publication.
So, let us take the proposition $B$\- (consequent of the implication) in the
form: $B=diag\left(1,0,0,0\right)$, that means that consequent is both true
and believable proposition. Then the expression Eq. (10) implies that matrix
$\left(A\Longrightarrow B\right)$ has the form:
$(A\Longrightarrow B)=diag\left(p_{2}+p_{4},p_{1}+p_{3},0,0\right),$ (11)
and hence its projections in DRS (1) and HRS (2) systems are respectively :
$\left(A\Longrightarrow B\right)_{1}=\begin{pmatrix}1&\\\ &0\end{pmatrix}$ ,
and $\left(A\Longrightarrow B\right)_{2}=\begin{pmatrix}p_{2}+p_{4}&\\\
&p_{1}+p_{3}\end{pmatrix}$
On the other hand if one choose the consequent $B$ in the form
$B=diag\left(0,1,0,0\right)$ that means that consiquent $B$ is true but
unbelievable proposition then according to expression $\left(10\right)$ one
obtain for the implication $\left(A\Longrightarrow B\right)$ the relation:
$\left(A\Longrightarrow B\right)=diag\left(0,1,0,0\right),$ (12)
and hence the projections of this proposition in two cognitive systems
are:$\left(A\Longrightarrow B\right)_{1}=\begin{pmatrix}1&\\\ &0\end{pmatrix}$
and $\left(A\Longrightarrow B\right)_{2}=\begin{pmatrix}0&\\\
&1\end{pmatrix}.$
Now if we make the natural assumption that after the first ( unconscious)
stage of reasoning, when two cognitive systems operate jointly, at the second
stage the conscious evaluation of the validity of a conclusion $V$ occurs in
accordance with the simple rule:
$V=aP_{t}+\left(1-a\right)P_{b},$ (13)
(where $a$ $(0\leqslant a\leqslant 1)$ is certain number coefficient depending
on age,intellect,training of the subject and possibly some other factors).Note
that this assumption in fact coincides with similar rule which was used in the
paper 4s . Now returning to the above example of interest we result in that
the magnitude of the bias-belief effect $V$ can be evaluate quantitatively as
$V\equiv V_{1}-V_{2}=\left(1-a\right)\left(p_{2}+p_{4}\right)$.We believe that
although the value of coefficient $a$ is unknown in advance nevertheless the
validity of the Eq. (13) can be explicitly verified in seria of properly
organized psychological experiments with various subjects using the identical
cognitive tasks .
In conclusion of our study let us formulate once more the central results of
the present paper:
1)We introduced the novel version of DCL with both discrete and continuous
logical operations between generalized propositions .
2)We proposed the concrete interpretation of propositions in DCL as integral
mental structures that include both logical and heuristic constituents.
3)We stated the specific uncertainty relation between logic rigour and
heuritic grasp that reflect complementary aspects of human reasoning process.
4)We proposed phenomenological model of human reasoning based on simplified
version of DCL and demonstrated that it is able to explain belief-bias effect
qualitatively and possibly quantitatively as well.
All these conclusions we hope to discuss more detail in our further
publications.
## References
* (1) Stent GS.Proc Am Philos Soc. 2004 Jun;148(2): 205-12.
* (2) N. Bohr, Science, New Series, Vol. 111, No. 2873 (Jan. 20, 1950), pp. 51-54
* (3) Enz, P.; von Meyenn, Karl (editors); Schlapp, Robert (translator) Wolfgang Pauli, Writings on physics and philosophy, pp218 - 279, Berlin: Springer Verlag,(1994)
* (4) Jonathan St. B. T. Evans, Thinking & Reasoning, V13, Issue 4, pp321-339, (2007)
* (5) E. D. Vol, physics.soc-ph 1306.2433, (2013)
|
arxiv-papers
| 2013-09-30T07:11:05 |
2024-09-04T02:49:51.690984
|
{
"license": "Public Domain",
"authors": "E. D. Vol",
"submitter": "Evgenii D. Vol",
"url": "https://arxiv.org/abs/1309.7740"
}
|
1309.7742
|
# Emergent of Majorana Fermion mode and Dirac Equation in Cavity Quantum
Electrodynamics
Sujit Sarkar Poornaprajna Institute of Scientific Research, 4 Sadashivanagar,
Bangalore 5600 80, India.
###### Abstract
We present the results of low lying excitation of coupled optical cavity
arrays. We derive the Dirac equation for this system and explain the existence
of Majorana fermion mode in the system. We present quite a few analytical
relations between the Rabi frequency oscillation and the atom-photon coupling
strength to achieve the different physical situation of our study and also the
condition for massless excitation in the system. We present several analytical
relations between the Dirac spinor field, order and disorder operators for our
systems. We also show that the Luttinger liquid physics is one of the
intrinsic concept in our system.
PACS: 42.50.Pq, 03.65.Vf, 42.50.-p
Keywords: Cavity Quantum Electrodynamics , Phases: Geometric, Dynamic or
Topological , Quantum Optics
Introduction: The recent experimental success in engineering strong
interaction between the photons and atoms in high quality micro-cavities opens
up the possibility to use light matter system as quantum simulators for many
body physics [1-21]. Many interesting results are coming out to understand the
complicated quantum many body system. The Bose-Hubbard model, quantum spin
model and the other exotic quantum phases of the quantum many body system have
already been studied [3-7]. The further application of the basic principle of
cavity-QED system is the circuit QED [22-26]. When a qubit (qubits) coupled to
the high quality LC circuits, it presents the same physical picture with few
extra achievement over the conventional cavity-QED system. A focus on the
coupled cavities is one of the most potential candidate for an efficient
quantum simulator due to the control of the microcavities parameters and
success of fabrication of large scale cavity arrays [25-26]. This is the very
brief discussion of the presence status of the cavity QED system.
In the present study one of our goal is to predict the presence of Majorana
fermions in our model system. Before we proceed further , we would like to
describe very briefly about the appreance of Majorana fermions in quantum
condensed matter system. Majorana had introduced a special kind of fermions
which are their own antiparticle, i.e., the neutral particle majo ; wil . He
had introduced this particle to describe neutrions. In recent years, there are
several candidates of Majorana fermions in quantum condensed matter system
like quantum Hall system with filling fraction $5/2$ read ; read2 . Kitaev at
first found the existence of Majorana fermion mode in one dimensional model
kitaev . Many research group have already been proposed the physically
existence of MFs at the edge state of 1D system like electrostatic defects
lines in superconductor, quasi-one dimensional superconductor and cold atom
trapped in one dimension wimmer ; kop . Majorana fermions are obey non-Abelian
statics both in 2D and 1D, allowing of certain gate operation required in
quantum computation nayak . Due to the non-local character, the qubit built
out of Majorana fermions are insensitive to local parity conserving
perturbation nayak ; lut ; oreg ; potter ; jiang . The search for
experimentally accessible systems that are described by a Dirac equation has
received much attention in recent years nayak ; lut ; oreg ; potter ; jiang .
In this research paper, we present an extensive derivation of Dirac equation
and also the existence of Majorana fermions mode in an optical cavity array.
We also present the analytical relation between the Rabi frequency oscillation
and the atom-photon coupling strengths to mimic the transverse Ising model,
Dirac equation, magnetic ordered state, quantum paramagnetic state and
massless excitation. Quantum state engineering of the optical cavity array
system is in the state of art due to the rapid technical development of this
field [1] therefore one can achieve these quantum phases in the laboratory.
The Model Hamiltonian and Majorana Fermion Modes:
The Hamiltonian of our present study consists of three parts:
$H~{}=~{}{H_{A}}~{}+~{}{H_{C}}~{}+~{}{H_{AC}}$ (1)
The Hamiltonians are the following
${H_{A}}~{}=~{}\sum_{j=1}^{N}{{\omega}_{e}}|e_{j}><e_{j}|~{}+~{}{\omega}_{ab}|b_{j}><b_{j}|$
(2)
where $j$ is the cavity index. ${\omega}_{ab}$ and ${\omega}_{e}$ are the
energies of the state $|b>$ and the excited state respectively. The energy
level of state $|a>$ is set as zero. $|a>$ and $|b>$ are the two stable state
of a atom in the cavity and $|e>$ is the excited state of that atom in the
same cavity. The following Hamiltonian describes the photons in the cavity,
${H_{C}}~{}=~{}{{\omega}_{C}}\sum_{j=1}^{N}{{a_{j}}}^{\dagger}{a_{j}}~{}+~{}{J_{C}}\sum_{j=1}^{N}({{a_{j}}}^{\dagger}{a_{j+1}}+h.c),$
(3)
where ${a_{j}}^{\dagger}({a_{j}})$ is the photon creation (annihilation)
operator for the photon field in the $j$’th cavity, ${\omega}_{C}$ is the
energy of photons and $J_{C}$ is the tunneling rate of photons between
neighboring cavities. The interaction between the atoms and photons and also
by the driving lasers are described by
${H_{AC}}~{}=~{}\sum_{j=1}^{N}[(\frac{{\Omega}_{a}}{2}e^{-i{{\omega}_{a}}t}+{g_{a}}{a_{j}})|e_{j}><a_{j}|+h.c]+[a\leftrightarrow
b].$ (4)
Here ${g_{a}}$ and ${g_{b}}$ are the couplings of the cavity mode for the
transition from the energy states $|a>$ and $|b>$ to the excited state.
${\Omega}_{a}$ and ${\Omega}_{b}$ are the Rabi frequencies of the lasers with
frequencies ${\omega}_{a}$ and ${\omega}_{b}$ respectively.
The authors of Ref. hart1 ; hart2 ; sujop have derived an effective spin
model by considering the following physical processes: A virtual process
regarding emission and absorption of photons between the two stable states of
neighboring cavity yields the resulting effective Hamiltonian as
${H_{xy}}=\sum_{j=1}^{N}B{{\sigma}_{j}}^{z}~{}+~{}\sum_{j=1}^{N}(\frac{J_{1}}{2}{{\sigma}_{j}}^{\dagger}{{\sigma}_{j+1}}^{-}~{}+~{}\frac{J_{2}}{2}{{\sigma}_{j}}^{-}{{\sigma}_{j+1}}^{-}+h.c)$
(5)
When $J_{2}$ is real then this Hamiltonian reduces to the XY model. Where
${{\sigma}_{j}}^{z}=|b_{j}><b_{j}|~{}-~{}|a_{j}><a_{j}|$,
${{\sigma}_{j}}^{+}=|b_{j}><a_{j}|$, ${{\sigma}_{j}}^{-}=|a_{j}><b_{j}|$ .
$\displaystyle H_{xy}$ $\displaystyle=$
$\displaystyle\sum_{i=1}^{N}(B{{\sigma}_{i}}^{z}~{}+~{}{J_{1}}({{\sigma}_{i}}^{x}{{\sigma}_{i+1}}^{x}+{{\sigma}_{i}}^{y}{{\sigma}_{i+1}}^{y})$
(6)
$\displaystyle+{J_{2}}({{\sigma}_{i}}^{x}{{\sigma}_{i+1}}^{x}-{{\sigma}_{i}}^{y}{{\sigma}_{i+1}}^{y}))$
$\displaystyle=\sum_{i=1}^{N}B({{\sigma}_{i}}^{z}~{}+~{}{J_{x}}{{\sigma}_{i}}^{x}{{\sigma}_{i+1}}^{x}~{}+~{}{J_{y}}{{\sigma}_{i}}^{y}{{\sigma}_{i+1}}^{y}).$
With ${J_{x}}=(J_{1}+J_{2})$ and ${J_{y}}=(J_{1}-J_{2})$.
We follow the references james ; hart1 , to present the analytical expression
for the different physical parameters of the system.
$B=\frac{\delta_{1}}{2}-\beta$ (7) $\displaystyle\beta$ $\displaystyle=$
$\displaystyle\frac{1}{2}[\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}({\Delta}_{b}-\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}-\frac{{|{\Omega_{b}}|}^{2}}{4({\Delta}_{a}-{\Delta}_{b})}-{\gamma_{b}}{g_{b}}^{2}-{\gamma_{1}}{g_{a}}^{2}$
(8)
$\displaystyle+{\gamma_{1}}^{2}\frac{{g_{a}}^{4}}{{\Delta_{b}}}-(a\leftrightarrow
b)]$
${J_{1}}=\frac{\gamma_{2}}{4}(\frac{{|{\Omega_{a}}|}^{2}{g_{b}}^{2}}{{{\Delta}_{a}}^{2}}+\frac{{|{\Omega_{b}}|}^{2}{g_{a}}^{2}}{{{\Delta}_{b}}^{2}}),{J_{2}}=\frac{\gamma_{2}}{2}(\frac{{\Omega_{a}}{\Omega_{b}}g_{a}g_{b}}{{\Delta}_{a}{\Delta_{b}}}).$
(9)
Where $\gamma_{a,b}=\frac{1}{N}\sum_{k}\frac{1}{{\omega}_{a,b}-{\omega}_{k}}$
$\gamma_{1}=\frac{1}{N}\sum_{k}\frac{1}{({\omega}_{a}+{\omega}_{b})/2-{\omega}_{k}}$
and
$\gamma_{2}=\frac{1}{N}\sum_{k}\frac{e^{ik}}{({\omega}_{a}+{\omega}_{b})/2-{\omega}_{k}}$
${\delta_{1}}={\omega}_{ab}-({\omega}_{a}-{\omega}_{b})/2$,
${\Delta}_{a}={\omega}_{e}-{\omega}_{a}$.
${\Delta}_{b}={\omega}_{e}-{\omega}_{a}-({\omega}_{ab}-{\delta_{1}})$.
${{\delta}_{a}}^{k}={\omega}_{e}-{\omega}_{k}$,
${{\delta}_{b}}^{k}={\omega}_{e}-{\omega}_{k}-({\omega}_{ab}-{\delta_{1}})$,
${\omega}_{k}={\omega}_{c}+J_{c}\sum_{k}cosk$. $g_{a}$ and $g_{b}$ are the
couplings of respective transition to the cavity mode, ${\Omega}_{a}$ and
${\Omega}_{b}$ are the Rabi frequency of laser with frequency $\omega_{a}$ and
$\omega_{b}$.
The system reduces to Ising model with transverse field at $J_{1}=J_{2}$,
i.e., $J_{x}$ become $J_{1}+J_{2}$ and $J_{y}=0$. The effective Hamiltonian
become the transverse Ising model which studied in the previous literature ss
; mussardo ; druff . Here our main motivation is to use some of important
results of this model Hamiltonian to discuss the relevant physics of array of
cavity QED system.
Before we proceed further, we would like to discuss in detail the analytical
relation between the different coupling constants of cavity QED system to
achieve this Hamiltonian. In the microcavity array, the condition for
$J_{1}=J_{2}$ achieve when
${{\Omega}_{a}}^{2}{g_{b}}^{2}{\Delta_{b}}^{2}+{{\Omega}_{b}}^{2}{g_{a}}^{2}{\Delta_{a}}^{2}=2{\Omega}_{a}{\Omega}_{b}g_{a}g_{b}{\Delta}_{a}{\Delta}_{b}.$
(10)
The above condition implies that
${\Omega}_{a}={\Omega}_{b}\frac{g_{a}\Delta_{a}}{g_{b}\Delta_{b}}$. The only
constraint is that ${\Delta}_{a}\neq{\Delta}_{b}$, the magnetic field diverge
when ${\Delta}_{a}={\Delta}_{b}$. At the same time,
${\Omega}_{a}={\Omega}_{b}$ and $g_{a}=g_{b}$ are also not possible because
this limit also leads to the condition ${\Delta}_{a}={\Delta_{b}}$. Suppose we
consider, ${\Omega}_{a}=\alpha_{1}{\Omega}_{b}$, $g_{a}=\alpha_{2}g_{b}$ and
${\Delta}_{a}=\alpha_{3}{\Delta}_{b}$. These relations implies that
${\alpha_{1}}^{2}+{\alpha_{2}}^{2}{\alpha_{3}}^{2}=2{\alpha}_{1}{\alpha}_{2}{\alpha}_{3}$.
${\alpha}_{1}={\alpha}_{2}{\alpha}_{3}$, ${\alpha}_{1},{\alpha_{2}}$ and
${\alpha_{3}}$ are the numbers. These analytical relations help to implement
the transverse Ising model Hamiltonian but $\alpha_{1}$, $\alpha_{2}$ and
$\alpha_{3}$ should not be equal to 1.
The quantum state engineering of cavity QED is in the state of art due to the
rapid progress of technological development of this field [1]. Therefore one
can achieve this limit to get the desire quantum state.
$H_{T}=B\sum_{j=1}^{N}({{\sigma}_{z}}(j)+\lambda{{\sigma}_{x}}(j){{\sigma}_{x}}(j+1)),$
(11)
where $\lambda=\frac{J_{1}+J_{2}}{B}$. The transverse Ising model was studied
widely in the literature and also exhibit a quantum phase transition between
the magnetically ordered state to the quantum paramagnetic phase for
$\lambda>1$ and $\lambda<1$ respectively ss ; mussardo ; druff .
Now we express the condition for the magnetic order phase and quantum
paramagnetic phase in terms of the physical parameters of the optical cavity
QED system which gives us the relevant physics of the system.
The condition for the magnetic ordered system can be expressed as
$\frac{\gamma_{2}}{4}(\frac{{\Omega}_{a}{g_{b}}^{2}}{{\Delta_{a}}^{2}}+\frac{{\Omega}_{b}{g_{a}}^{2}}{{\Delta_{b}}^{2}})+\frac{\gamma_{2}}{2}(\frac{\Omega_{a}\Omega_{b}g_{a}g_{b}}{\Delta_{a}\Delta_{b}})>\omega_{ab}-\frac{\omega_{a}-\omega_{b}}{2}-2\beta$
(12)
The condition for the quantum paramagnetic phase is
$\frac{\gamma_{2}}{4}(\frac{{\Omega}_{a}{g_{b}}^{2}}{{\Delta_{a}}^{2}}+\frac{{\Omega}_{b}{g_{a}}^{2}}{{\Delta_{b}}^{2}})+\frac{\gamma_{2}}{2}(\frac{\Omega_{a}\Omega_{b}g_{a}g_{b}}{\Delta_{a}\Delta_{b}})<\omega_{ab}-\frac{\omega_{a}-\omega_{b}}{2}-2\beta$
(13)
When the applied magnetic field is absent, the effective Ising model has two
degenerate ground states. The ground states are
$|A>=\Pi_{j}{|\rightarrow>}_{j}$, $|B>=\Pi_{j}{|\leftarrow>}_{j}$. For a
finite magnetic field but less than $J_{1}+J_{2}$, the system has a tendency
to flip the pseudo spin. At that phase one can write down the true eigen
state, $|{\psi}_{A}>=\frac{1}{\sqrt{2}}(|A>+|B>)$ ,
$|{\psi}_{B}>=\frac{1}{\sqrt{2}}(|A>-|B>)$. Now our main intention is to
recast this spin model in spinless fermion model through Jordon-Wigner
transformation which relate the spin operators to the spinless fermion
operators. We use the following relation:
${\sigma}_{z}=2c^{\dagger}(j)c(j)-1$,
${\sigma}_{x}(j){\sigma}_{x}(j+1)=(c^{\dagger}(n)-c(n))(c^{\dagger}(n+1)-c(n+1))$.
One can write the Hamiltonian after the Jordon-Wigner transformation as
$H=2\sum_{j=1}^{N}{c}^{\dagger}(j)c(j)+\lambda(c^{\dagger}(j)-c(j))(c^{\dagger}(j+1)-c(j+1))$
(14)
We solve this Hamiltonian, to get the energy spectrum by taking the Fourier
transform.
$c(j)=\frac{1}{\sqrt{N}}\sum_{k}e^{-ika}$,
$c^{\dagger}(j)=\frac{1}{\sqrt{N}}\sum_{k}e^{ika}$. Where $c_{k}$ and
${c_{k}}^{\dagger}$ are the fermionic annihilation and creation operator in
momentum space.
The Hamiltonian reduce to
$\displaystyle H$ $\displaystyle=$ $\displaystyle 2\sum_{k>0}(1+\lambda
cosk)({c_{k}}^{\dagger}{c_{k}}+{c_{-k}}^{\dagger}{c_{-k}})$ (15)
$\displaystyle+2i\lambda\sum_{k>0}sink({c_{k}}^{\dagger}{c_{-k}}^{\dagger}+{c_{k}}{c_{-k}})$
Now our main task is to express the Hamiltonian in the diagonalized form. We
follow the Bogoliubov transformation.
$\eta_{k}=\alpha_{k}c_{k}+i{\beta}_{k}{c_{-k}}^{\dagger}$ and
$\eta_{-k}=\alpha_{k}c_{-k}-i{\beta}_{k}{c_{k}}^{\dagger}$, $k>0$.
The operator ${\eta}_{k}$ and ${{\eta}_{k}}^{\dagger}$ are the fermionic
operators. We use the following relations,
$\\{{\eta}_{k},{{\eta}_{p}}^{\dagger}\\}=\delta_{k,p}$,
$\\{{\eta}_{k},{{\eta}_{p}}\\}=0$,
$\\{{{\eta}_{k}}^{\dagger},{{\eta}_{p}}^{\dagger}\\}=0$. This relation
implies, ${{\alpha}_{k}}^{2}+{{\beta}_{k}}^{2}=1$. One can also revert the
relation between $c_{k}$ and ${\eta}_{k}$. We also parameterize
$\alpha_{k}=cos\theta_{k}$ and $\beta_{k}=sin\theta_{k}$. One can express the
transformed Hamiltonian in two parts
$H=H_{A}+H_{B}$ (16) $\displaystyle H_{A}$ $\displaystyle=$
$\displaystyle\sum_{k>0}[-2(1+\lambda
cosk)({{\alpha}_{k}}^{2}-{{\beta}_{k}}^{2})+4\lambda
sink{\alpha}_{k}{\beta}_{k}]$ (17)
$\displaystyle({{\eta}_{k}}^{\dagger}{\eta_{k}}{{\eta}_{-k}}^{\dagger}{\eta_{-k}})$
$\displaystyle H_{B}$ $\displaystyle=$ $\displaystyle\sum_{k>0}[4i(1+\lambda
cosk){\alpha_{k}}{\beta_{k}}+2i\lambda
sink({{\alpha}_{k}}^{2}-{\beta_{k}}^{2})]$ (18)
$\displaystyle({{\eta}_{k}}^{\dagger}{\eta_{-k}}^{\dagger}{{\eta}_{k}}{\eta_{-k}})$
To express this Hamiltonian in the diagonal form, we find the following
relation
$4(B+\frac{\gamma_{2}{\Omega_{b}}^{2}{g_{a}}^{2}}{{\Delta_{b}}^{2}}cosk){\alpha_{k}}{\beta_{k}}+2\frac{\gamma_{2}{\Omega_{b}}^{2}{g_{a}}^{2}}{{\Delta_{b}}^{2}}sink({\alpha_{k}}^{2}-{\beta_{k}}^{2})=0$.
Finally this gives the condition,
$tan2{\theta_{k}}=\frac{2\gamma_{2}{\Omega_{b}}^{2}{g_{a}}^{2}}{2{\gamma_{2}}{\Omega_{b}}^{2}{g_{a}}^{2}cosk-({\delta_{1}}-2\beta){\Delta_{b}}^{2}}$
(19)
$2{\alpha}_{k}{\beta}_{k}=sin2{\theta_{k}}$,
${\alpha_{k}}^{2}-{\beta_{k}}^{2}=cos2\theta_{k}$. Now we analysis the
spectrum:
$H=2\sum_{k}\Omega_{K}{\eta_{k}}^{\dagger}{\eta_{k}}$ (20)
We can express this energy spectrum in terms of Rabi frequency oscillation,
atom-photon coupling strength,
$\Omega_{k}=\frac{2}{{\delta_{1}}/2-\beta}\sqrt{{({\delta_{1}}/2-\beta)}^{2}+{\gamma_{2}}^{2}\frac{{{\Omega_{b}}}^{4}{g_{a}}^{4}}{{\Delta_{a}}^{4}}+2({\delta_{1}}/2-\beta){\gamma_{2}}\frac{{{\Omega_{b}}}^{2}{g_{a}}^{2}}{{\Delta_{b}}^{2}}}$
(21)
The minimum occurs at $k=\pm\pi$,
$\Omega_{k=\pm\pi}=2|1-\lambda|=2|1-\frac{2\gamma_{2}}{\delta_{1}-2\beta}\frac{{{\Omega}_{b}}^{2}{g_{a}}^{2}}{{\Delta_{b}}^{2}}|$.
We are interested in the continuum limit and also restore the lattice spacing
$\alpha$ and measure the momentum w.r.t the minimum value,
$k=\pi+k^{{}^{\prime}}\alpha$ . The energy expression which contains the
physical dimension of energy is
$E(k^{{}^{\prime}})=\frac{\Omega_{k}}{2\alpha}$. In this limit,
$E(k^{{}^{\prime}})=\sqrt{{(\frac{1-\lambda}{\alpha})}^{2}+\lambda{k^{{}^{\prime}}}^{2}}$.
If $\lambda$ is close to a critical value $\lambda\sim 1$, we then have the
dispersion of a particle with mass $m=\frac{1-\lambda}{\alpha}$.
If $\lambda=1$, it becomes the massless particle $E(k^{{}^{\prime}})\sim
k^{{}^{\prime}}$.
In the cavity QED system, we can express the condition for massless excitation
of the system as
$({{\delta_{1}}/2-\beta}){\Delta_{b}}^{2}={\gamma_{2}}{\Omega_{b}}^{2}{g_{a}}^{2}$.
In terms of fields, $\eta(a)=\frac{1}{\sqrt{N}}\sum_{k}e^{ika}{\eta}_{k}$.
${\eta(a)}^{\dagger}=\frac{1}{\sqrt{N}}\sum_{k}e^{-ika}{{\eta}_{k}}^{\dagger}$.
${\chi}_{1}(a)=(1/2)(\eta(a)+{\eta(a)}^{\dagger})$,
${\chi}_{2}(a)=(1/2i)(\eta(a)-{\eta(a)}^{\dagger})$
${{\chi}_{1}(a)}^{\dagger}={\chi}_{1}(a)$,
${{\chi}_{2}(a)}^{\dagger}={\chi}_{2}(a)$.
$\\{{\chi_{1}}(x1),{\chi_{2}}(x2)\\}={\delta}_{x1,x2}{\delta}_{1,2}$.
Therefore, ${\chi}_{1}$ and ${\chi}_{2}$ are satisfying the all properties of
neutral fermionic fields what Majorana proposed.
The authors of choi2 have investigated the low lying excitation of one
dimensional array of circuit QED (cktQED) with each cktQED being in the ultra
strong coupling regime and they have found the Majorana bound state. But the
starting Hamiltonian of our system is completely different.
Now we calculate the energy density of the system using Eq. 21. We would like
to integrate the dispersion spectrum ${\Omega}_{k}$ to get the energy density.
The analytical expression for energy density is
${\epsilon}_{0}=\frac{2}{\pi}~{}(1+\frac{2\gamma_{2}}{\delta_{1}-2\beta}\frac{{{\Omega}_{b}}^{2}{g_{a}}^{2}}{{\Delta_{b}}^{2}})~{}E(\frac{\pi}{2},\sqrt{1-{\gamma}^{2}})$
(22)
Where
${\gamma}=|\frac{(\delta_{1}-2\beta){\Delta_{b}}^{2}-2{\gamma_{2}}{\Omega_{b}}^{2}{g_{a}}^{2}}{(\delta_{1}-2\beta){\Delta_{b}}^{2}+2{\gamma_{2}}{\Omega_{b}}^{2}{g_{a}}^{2}}|$,
$E(\frac{\pi}{2},\sqrt{1-{\gamma}^{2}})$ is the complete Elliptic integral of
2nd kind. After a little bit of calculation, we obtain the analytical
expression in the asympototic limit for energy density.
$\displaystyle{\epsilon}_{0}$ $\displaystyle=$ $\displaystyle
1+1/2~{}(ln|4\frac{(\delta_{1}-2\beta){\Delta_{b}}^{2}+2{\gamma_{2}}{\Omega_{b}}^{2}{g_{a}}^{2}}{(\delta_{1}-2\beta){\Delta_{b}}^{2}-2{\gamma_{2}}{\Omega_{b}}^{2}{g_{a}}^{2}}|-1/2){\gamma}^{2}$
(23)
$\displaystyle+3/16(ln|4\frac{(\delta_{1}-2\beta){\Delta_{b}}^{2}+2{\gamma_{2}}{\Omega_{b}}^{2}{g_{a}}^{2}}{(\delta_{1}-2\beta){\Delta_{b}}^{2}-2{\gamma_{2}}{\Omega_{b}}^{2}{g_{a}}^{2}}|-13/12){\gamma}^{4}$
Derivation of Dirac Equation and Condition of Massless Excitations:
During the derivation of Dirac equation we rotate the y-axis of the spin basis
through $\pi/2$. Through this rotation z-axis become x-axis and x-axis become
-z-axis. The physics of the system remain the same. The starting Hamiltonian
of our system now becomes
$H=\sum_{n}[\lambda s_{z}(n)s_{z}(n+1)+s_{x}(n)].$ (24)
We recast the Hamiltonian in the following form because it will help us to use
the order and disorder operator directly to the derivation of the equation of
motion and finally the Dirac equation. Here we present an extensive derivation
of Dirac equation for our model system.
Kramers-Wannier symmetry for two dimensional Ising model reveals for the case
of one-dimensional quantum Ising chain through the dual lattice to the
original one (site index n). Here we introduce the order and disorder operator
following mussardo ; druff . These operators are defining the sites of the
dual lattice, i.e., we define the operator between the nearest-neighbor site
of the original lattice. The analytical relation between the Pauli operators
and $\mu$ operators are the following:
${{\mu}_{z}}^{2}=1={{\mu}_{x}}^{2},$ (25)
${{\mu}_{z}}(n-1/2){{\mu}_{z}}(n+1/2)={\sigma}_{x}(n).$ (26)
${{\mu}_{x}}(n+1/2)={\sigma}_{z}(n){\sigma}_{z}(n+1),$ (27)
${{\mu}_{z}}(n+1/2)=\Pi_{j=1}^{n}{\sigma_{x}}(j).$ (28)
${{\sigma_{z}}}(n)=\Pi_{j=0}^{n-1}{\mu_{x}}(j+1/2),$ (29)
$[{\mu}_{x}(n+1/2),{\mu}_{z}(n^{{}^{\prime}}+1/2)]=2\delta_{n,n^{{}^{\prime}}}$
(30) $[{\mu}_{z}(n+1/2),{\mu}_{z}(n^{{}^{\prime}}+1/2)]=0,$ (31)
$[{\mu}_{z}(n+1/2),{\sigma}_{x}(n^{{}^{\prime}})]=0$ (32)
The operator ${\mu}_{z}(n+1/2)$ acting on the original spin of the lattice
makes a spin flip of all those spin placed on the left hand side of spin at
the site n. Therefore ${\mu}_{z}(n+1/2)$ is a kink operator, it introduce the
disorder in the system. It is very clear from the above analytical relation of
the operators that ${\mu}_{x}$ is related with the allinegment of the spin
operator.
Here we define the Dirac spinor, $\chi_{1}(n)=\sigma_{z}(n)\mu_{z}(n+1/2)$ and
$\chi_{2}(n)=\sigma_{z}(n)\mu_{z}(n-1/2)$.
Now our main task is to find the equation of motion for the operators,
${\sigma}_{3}(n)$ and ${\mu}_{3}(n)$ which help us to build the Dirac
equation.
The equation of motion for the ${\sigma}_{z}(n)$ is the following:
$\frac{\partial\sigma_{z}(n)}{\partial\tau}=[H,\sigma_{z}(n)]={\sigma_{x}(n)\sigma_{z}(n)}$
(33)
The equation of motion for $\mu_{z}(n+1/2)$ is the following:
$\displaystyle\frac{\partial\mu_{z}(n+1/2)}{\partial\tau}$ $\displaystyle=$
$\displaystyle\lambda{\mu_{x}(n+1/2)\mu_{z}(n+1/2)}$ (34)
$\displaystyle=\lambda\sigma_{z}(n)\sigma_{z}(n+1/2)\mu_{z}(n+1/2)$
Now we use the properties of the $\sigma$ and $\mu$ operators to derive the
equation of motion for the Majorana fields $\chi_{1}(n)$ and $\chi_{2}(n)$.
$\frac{\partial\chi_{1}(n)}{d\tau}=\frac{\partial\sigma_{z}(n)}{\partial\tau}\mu_{z}(n+1/2)+\sigma_{z}(n)\frac{\partial\mu_{z}(n)}{\partial\tau}.$
(36)
$\frac{\partial\chi_{1}(n)}{d\tau}=\sigma_{x}(n)\sigma_{z}(n)\mu_{z}(n+1/2)+\lambda\sigma_{z}(n)\sigma_{z}(n)\sigma_{z}(n+1)\mu_{z}(n+1/2).$
(37) $\displaystyle\frac{\partial\chi_{1}(n)}{d\tau}$ $\displaystyle=$
$\displaystyle-\sigma_{z}(n)\mu_{z}(n-1/2)\mu_{z}(n+1/2)\mu_{z}(n+1/2)$ (38)
$\displaystyle+\lambda\sigma_{z}(n)\sigma_{z}(n)\sigma_{z}(n+1)\mu_{z}(n+1/2).$
$\frac{\partial\chi_{1}(n)}{d\tau}=-\chi_{2}(n)+\lambda\chi_{2}(n+1).$ (39)
Now the equations of motion for $\chi_{2}(n)$ are
$\displaystyle\frac{\partial\chi_{2}(n)}{d\tau}$ $\displaystyle=$
$\displaystyle\frac{\partial\sigma_{z}(n)}{\partial\tau}\mu_{z}(n-1/2)$ (40)
$\displaystyle+\sigma_{z}(n)\frac{\partial\mu_{z}(n-1/2)}{\partial\tau}$
$\displaystyle\frac{\partial\chi_{2}(n)}{d\tau}$ $\displaystyle=$
$\displaystyle\sigma_{x}(n)\sigma_{z}(n)\mu_{z}(n-1/2)$ (41)
$\displaystyle+\lambda\sigma_{z}(n)\sigma_{z}(n-1)\sigma_{z}(n)\mu_{z}(n-1/2)$
$\displaystyle\frac{\partial\chi_{2}(n)}{d\tau}$ $\displaystyle=$
$\displaystyle\mu_{z}(n-1/2)\mu_{z}(n+1/2)\sigma_{z}(n)\mu_{z}(n-1/2)$ (42)
$\displaystyle+\lambda\sigma_{z}(n-1)\mu_{z}(n-1/2).$
After a little bit of calculations and using the relation between the disorder
operators (Eq. 23-30), we finally arrive the equation of motion of
$\chi_{2}(n)$ as,
$\frac{\partial\chi_{2}(n)}{d\tau}=-\chi_{1}(n)+\lambda\chi_{1}(n-1).$ (43)
These two fields, $\chi_{1}(n)$ and $\chi_{2}(n)$ satisfy the following
relations, $\\{\chi_{1}(n1),\chi_{2}(n2)\\}=2\delta_{n1,n2}$. One can write
down the above equation in the following compact form,
$({\gamma}^{0}\frac{\partial}{\partial t}+{\gamma}^{3}\frac{\partial}{\partial
r}+m)\chi=0.$ (44)
where ${\chi}^{\dagger}=(\chi_{1},\chi_{2})$ and $m=\frac{1-\lambda}{\alpha}$,
${\gamma}^{0}=\left(\begin{array}[]{cc}0&1\\\
1&0\end{array}\right),{\gamma}^{3}=\left(\begin{array}[]{cc}1&0\\\
0&-1\end{array}\right)$ .
Now we prove the presence of Luttinger liquid physics is intrinsic to the
optical cavity array system by comparing the analytical relation of Majorana
fermion operators with the order and disorder operator with the free Dirac
field in Abelian bosonization theory. Here we express the analytical relation
between the Majorana operators and the disorder and Pauli operators:
${\chi}_{1}(n)={\sigma}_{z}(n){\mu}_{z}(n+1/2)=-{\mu}_{z}(n+1/2){\sigma}_{z}(n)$
(45)
${\chi}_{2}(n)={\sigma}_{z}(n){\mu}_{z}(n-1/2)={\mu}_{z}(n-1/2){\sigma}_{z}(n)$
(46)
${\sigma}_{z}(n){\chi}_{2}(n)={\mu}_{z}(n-1/2)={\chi}_{2}(n){\sigma}_{z}(n)$
(47)
${\sigma}_{z}(n){\chi}_{1}(n)={\mu}_{z}(n+1/2)=-{\chi}_{1}(n){\sigma}_{z}(n)$
(48)
${\sigma}_{z}(n)={\mu}_{z}(n-1/2){\chi}_{2}(n)={\chi}_{2}(n){\mu}_{z}(n-1/2)$
(49)
${\mu}_{z}(n+1/2){\chi}_{1}(n)=-{\chi}_{1}(n){\mu}_{z}(n+1/2)={\sigma}_{z}(n)$
(50)
The above relations can be extended to account for arbitrary space separation
between different operators. Then one obtains the following sets of
commutation relations.
${\sigma}_{z}(x1){\mu}_{z}(x2)={\mu}_{z}(x2){\sigma}_{z}(x1)sign(x1-x2)$ (51)
${\sigma}_{z}(x1){\chi}(x2)={\chi}(x2){\sigma}_{z}(x1)sign(x1-x2)$ (52)
${\mu}_{z}(x1){\chi}(x2)=-{\chi}(x2){\sigma}_{z}(x1)sign(x1-x2)$ (53)
It is very clear from the above analytical relations that
${{\chi}_{1}}^{\dagger}={\chi}_{1}$ and ${{\chi}_{2}}^{\dagger}={{\chi}_{2}}$.
The above relation has similarity with the free Dirac field in Abalian
bosonization theory, where Dirac field operator is a local product of two
phase exponential depending on the scalar field and its dual gia ; gogo , as
one study the Luttinger liquid physics in Abelian bosonization theory.
Therefore the Luttinger liquid physics is the intrinsic to the optical
microcavity system.
Conclusions
We have presented an extensive derivation of Dirac equation and the existence
of Majorana fermion modes for the optical cavity array with the relation
between Rabi frequency oscillation and the atom photon coupling strength. We
have presented the condition for massless excitation. We have also presented
several analytical relations between the Majorana field, order and disorder
operator.
Acknowledgement: The author would like to acknowledge the discussions with
Prof. S. Girvin during the international workshop/school on Dirac Materials
and Chandrashekar lecture at ICTS (December, 2012) and the library of Raman
Research Institute (Mr. Manjunath). Finally, the author would like to thank
Dr. P. K. Mukherjee for reading the manuscript carefully.
## References
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|
arxiv-papers
| 2013-09-30T07:28:42 |
2024-09-04T02:49:51.698003
|
{
"license": "Public Domain",
"authors": "Sujit Sarkar",
"submitter": "Sujit Sarkar",
"url": "https://arxiv.org/abs/1309.7742"
}
|
1309.7777
|
# Quantum phase transition of light in coupled optical cavity arrays: A
renormalization group study
Sujit Sarkar Poornaprajna Institute of Scientific Research, 4 Sadashivanagar,
Bangalore 5600 80, India.
###### Abstract
We study the quantum phase transition of light of a system when atom trapped
in microcavities and interacting through the exchange of virtual photons. We
predict the quantum phase transition between the photonic Coulomb blocked
induce insulating phase and anisotropic exchange induced photonic superfluid
phase in the system due to the existence of two Rabi frequency oscillations.
The renormalization group equation shows explicitly that for this system there
is no self-duality. The system also shows two Berezinskii-Kosterlitz-Thouless
(BKT) transitions for the different physical situation of the system. The
presence of single Rabi frequency oscillation in the system leads to the BKT
transition where system shows the quantum phase transition from photonic
metallic state to the Coulomb blocked induced insulating phase. For the other
BKT transition when the z-component of exchange interaction is absent, the
system shows the transition from the photonic metallic state to the photonic
superfluid phase. We also predict the commensurate to incommensurate
transition under the laser field detuning.
Pacs: 42.50.Dv, 42.50.Pq, 03.67.Bg, 75.10.Jm
Keywords: Cavity Quantum Electrodynamics, Renormalization Group Theory,
Quantum Optics and Quantum Spin Model
Introduction:
The recent experimental success in engineering strong interaction between the
photons and atoms in high quality micro-cavities opens up the possibility to
use light matter system as quantum simulators for many body physics [1-18].
The authors of Ref. (hart1 ,hart2 , suj2 ) have shown that effective spin
lattice can be generated with individual atom in the microcavities that are
coupled to each other via exchange of virtual photons. The two states of spin
polarization are represented by the two long lived atomic levels in the
system.
A Many body Hamiltonians can be created and probed in coupled cavity arrays.
In our previous study, we have explained explicitly the basic physics of the
formation of micro-optical cavity suj2 . Atoms in the cavity are used for
detection and also for generation of interaction between photons in the same
cavity. This artificial system can act as a quantum simulator [1-6]. This
micro-cavity system shows the different quantum phases and quantum phase
transition (QPT) of photonic states of the system.
QPT associate with the fundamental changes that occurs in the macroscopic
nature of the matter at zero temperature due to the variation of some external
parameter. Quantum phase transitions are characterized by the drastic change
in the ground state properties of the system driven by the quantum
fluctuations subir ; indrani .
In this paper, we derive the renormalization group equation for the continuum
field theoretical study of the model Hamiltonian. We predict with the
existence of photonic metallic state, the anisotropic exchange induces the
photonic superfluid phase and photon blocked induces the insulating phase. The
transition from the photonic superfluid phase to the photon blocked induced
insulating phase is the second order quantum phase transition. Whereas the
transition from photonic metallic state to the photonic insulating state is
the Berezinski-Kosterlitz-Thouless (BKT) transition subir ; gia ; druf where
the anisotropic exchange interaction is absent. The other BKT transition
occurs when we consider the transition from photonic metallic state to the
photonic superfluid state when the z-component interaction is absent.
The quantum state of engineering of the cavity QED system is in the state of
art due to the rapid technological development in this field [1]. Therefore
one can achieve the all states through the proper manipulation of laser
frequencies, detuning field, Rabi frequency oscillations. In our study, the
photonic metallic state is the Luttinger liquid phase of the microcavities
array system. The photonic superfluid phase is originated from the anisotropic
exchange interactions which discuss in Ref.jor . To the best of our knowledge,
the explicit study of the different quantum phases through the RG equations
and the proposal of photonic metallic state is the first in the literature for
the Cavity QED system. Now we discuss in detail about the photon blocked
induce insulating phase. The photon blocked is a dressed state of strongly
coupled atom-cavity system. In this phase the inhivitation of resonant
absorption of a second photon if one photon is already resonantly absorbed.
This phenomena was first observed by Birnbaum $et~{}al.$ birn . In this
studies they have predicted the photon blocked in the single atom in the
cavity birn . Here we predict this phenomena for the array of cavities, where
the atoms in different cavities interacting through virtual photon
emission/absorption process.
We will see after the continuum field theoretical study that our model
Hamiltonian contains two strongly relevant and mutually nonlocal perturbation
over the Gaussian (critical) theory. In such a situation the strong coupling
fixed point is usually determined by the most relevant perturbation whose
amplitude grows up according to its Gaussian scaling dimensions and it is not
much affected by the less relevant coupling terms. However, this is not the
general rule if the two operators exclude each other, i.e., if the field
configurations which minimize one perturbation term do not minimize the other.
In this case interplay between the two competing relevant operators can
produce a novel quantum phase transition through a critical point or a
critical line. Therefore, we would like to study the RG equation to interpret
the quantum phases of the system.
Model Setup and Motivation
We would like to solve the quantum phases and phase boundaries of coupled
cavity arrays by establishing the quantum spin chain model of the system hart1
; hart2 ; suj2 . At first we discuss the relation between the spin operators
and the atom-photon system. Our starting point is the Janes-Cummings
Hamiltonian, $H~{}=g({\sigma}a^{\dagger}+h.c)$. Cavity mode represent by the
bosonic operators ($a,{a}^{\dagger}$ ) and atomic mode represent by the
fermionic operator (${\sigma},{\sigma}^{\dagger}$). Where $\sigma$ is the
Pauli matrix which transform one excitation from the radiation field to the
atomic field. Therefore, one can write
${\sigma}^{\dagger}=\left(\begin{array}[]{cc}0&1\\\ 0&0\end{array}\right)$
, similarly one can write for ${\sigma}^{-}$ with $<1|=(1,0)$ and $<0|=(0,1)$.
The excitation in this system only transfer between atom and photon in the
cavity. For a fixed number ($n$) of total excitation, one can express the
manifold: ${H_{n}}=\\{|0,n>,|1,n-1>\\}$ provided $n\geq 1$. Here $|0,n>$ and
$|1,n-1>$ are represent atom in the ground state with $n$ photon and excited
states of the atom with ($n-1$) photon respectively. We consider the initial
state $|e,n-1>$, we obtain the state $|g,n>$ by the following operation,
$|g,n>={\sigma}a^{+}|e,n-1>$. Therefore, we may write the following relation
based on the conservation of the number of excitation.
${\sigma}^{\dagger}a|e,n-1>=~{}0~{}={\sigma}a^{\dagger}|g,n>$ and
${\sigma}^{\dagger}a|g,n>=|e,n-1>$. These relations are nothing but the
properties of spin operators acting on the spinors in the z basis.
Micro-cavities of a photonic crystal are coupled through the exchange of
photons. Each cavity consists of one atom with three levels in the energy
spectrum, two of them are long lived and represent two spin states of the
system and the other represent excited states (fig. 1 and fig. 2 of Ref. hart1
; hart2 ). Externally applied laser and cavity modes couple to each atom of
the cavity. It may induce the Raman transition between these two long lived
energy levels. Under a suitable detuning between the laser and the cavity
modes, virtual photons are created in the cavity which mediate interactions
with another atom in a neighboring cavity. One can eliminate the excited
states by choosing the appropriate detuning between the applied laser and
cavity modes. Then one can achieve only two states per atom in the long lived
state and the system can be described by a spin-1/2 Hamiltonian hart1 ; hart2
.
The Hamiltonian of the system consists of three parts:
$H~{}=~{}{H_{A}}~{}+~{}{H_{C}}~{}+~{}{H_{AC}}$ (1)
Hamiltonians are the following
${H_{A}}~{}=~{}\sum_{j=1}^{N}{{\omega}_{e}}|e_{j}><e_{j}|~{}+~{}{\omega}_{ab}|b_{j}><b_{j}|$
(2)
where $j$ is the cavity index. ${\omega}_{ab}$ and ${\omega}_{e}$ are the
energies of the state $|b>$ and the excited state respectively. The energy
level of state $|a>$ is set as zero. $|a>$ and $|b>$ are the two stable state
of a atom in the cavity and $|e>$ is the excited state of that atom in the
same cavity. The following Hamiltonian describes photons in the cavity,
${H_{C}}~{}=~{}{{\omega}_{C}}\sum_{j=1}^{N}{{a_{j}}}^{\dagger}{a_{j}}~{}+~{}{J_{C}}\sum_{j=1}^{N}({{a_{j}}}^{\dagger}{a_{j+1}}+h.c),$
(3)
where ${a_{j}}^{\dagger}({a_{j}})$ is the photon creation (annihilation)
operator for the photon field in the $j$’th cavity, ${\omega}_{C}$ is the
energy of photons and $J_{C}$ is the tunneling rate of photons between
neighboring cavities. Interaction between the atoms and photons and also by
the driving lasers are described by
${H_{AC}}~{}=~{}\sum_{j=1}^{N}[(\frac{{\Omega}_{a}}{2}e^{-i{{\omega}_{a}}t}+{g_{a}}{a_{j}})|e_{j}><a_{j}|+h.c]+[a\leftrightarrow
b].$ (4)
Here ${g_{a}}$ and ${g_{b}}$ are the couplings of the cavity mode for the
transition from the energy states $|a>$ and $|b>$ to the excited state.
${\Omega}_{a}$ and ${\Omega}_{b}$ are the Rabi frequencies of the lasers with
frequencies ${\omega}_{a}$ and ${\omega}_{b}$ respectively.
The authors of Ref. hart1 ; hart2 have derived an effective spin model by
considering the following physical processes: A virtual process regarding
emission and absorption of photons between the two stable states of neghbior
cavity yields the resulting effective Hamiltonian as
${H_{xy}}=\sum_{j=1}^{N}B{{\sigma}_{j}}^{z}~{}+~{}\sum_{j=1}^{N}(\frac{J_{1}}{2}{{\sigma}_{j}}^{\dagger}{{\sigma}_{j+1}}^{-}~{}+~{}\frac{J_{2}}{2}{{\sigma}_{j}}^{-}{{\sigma}_{j+1}}^{-}+h.c)$
(5)
When $J_{2}$ is real then this Hamiltonian reduces to the XY model. Where
${{\sigma}_{j}}^{z}=|b_{j}><b_{j}|~{}-~{}|a_{j}><a_{j}|$,
${{\sigma}_{j}}^{+}=|b_{j}><a_{j}|$, ${{\sigma}_{j}}^{-}=|a_{j}><b_{j}|$
$\displaystyle H_{xy}$ $\displaystyle=$
$\displaystyle\sum_{i=1}^{N}(B{{\sigma}_{i}}^{z}~{}+~{}{J_{1}}({{\sigma}_{i}}^{x}{{\sigma}_{i+1}}^{x}+{{\sigma}_{i}}^{y}{{\sigma}_{i+1}}^{y})$
(6)
$\displaystyle+{J_{2}}({{\sigma}_{i}}^{x}{{\sigma}_{i+1}}^{x}-{{\sigma}_{i}}^{y}{{\sigma}_{i+1}}^{y}))$
$\displaystyle=\sum_{i=1}^{N}B({{\sigma}_{i}}^{z}~{}+~{}{J_{x}}{{\sigma}_{i}}^{x}{{\sigma}_{i+1}}^{x}~{}+~{}{J_{y}}{{\sigma}_{i}}^{y}{{\sigma}_{i+1}}^{y}).$
with ${J_{x}}=(J_{1}+J_{2})$ and ${J_{y}}=(J_{1}-J_{2})$.
Here we present the analytical expression of $B$, $J_{1}$ and $J_{2}$ in terms
of different physical parameters of the system. We follow the references hart1
; james .
$B=\frac{\delta_{1}}{2}-\beta$ (7) $\displaystyle\beta$ $\displaystyle=$
$\displaystyle\frac{1}{2}[\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}({\Delta}_{b}-\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}-$
(8)
$\displaystyle\frac{{|{\Omega_{b}}|}^{2}}{4({\Delta}_{a}-{\Delta}_{b})}-{\gamma_{b}}{g_{b}}^{2}-{\gamma_{1}}{g_{a}}^{2}+{\gamma_{1}}^{2}\frac{{g_{a}}^{4}}{{\Delta_{b}}}-(a\leftrightarrow
b)]$
${J_{1}}=\frac{\gamma_{2}}{4}(\frac{{|{\Omega_{a}}|}^{2}{g_{b}}^{2}}{{{\Delta}_{a}}^{2}}+\frac{{|{\Omega_{b}}|}^{2}{g_{a}}^{2}}{{{\Delta}_{b}}^{2}})$
(9)
${J_{2}}=\frac{\gamma_{2}}{2}(\frac{{\Omega_{a}}{\Omega_{b}}g_{a}g_{b}}{{\Delta}_{a}{\Delta_{b}}}).$
(10)
where $\gamma_{a,b}=\frac{1}{N}\sum_{k}\frac{1}{{\omega}_{a,b}-{\omega}_{k}}$
$\gamma_{1}=\frac{1}{N}\sum_{k}\frac{1}{({\omega}_{a}+{\omega}_{b})/2-{\omega}_{k}}$
and
$\gamma_{2}=\frac{1}{N}\sum_{k}\frac{e^{ik}}{({\omega}_{a}+{\omega}_{b})/2-{\omega}_{k}}$
${\delta_{1}}={\omega}_{ab}-({\omega}_{a}-{\omega}_{b})/2$,
${\Delta}_{a}={\omega}_{e}-{\omega}_{a}$.
${\Delta}_{b}={\omega}_{e}-{\omega}_{a}-({\omega}_{ab}-{\delta_{1}})$.
${{\delta}_{a}}^{k}={\omega}_{e}-{\omega}_{k}$,
${{\delta}_{b}}^{k}={\omega}_{e}-{\omega}_{k}-({\omega}_{ab}-{\delta_{1}})$,
$g_{a}$ and $g_{b}$ are the couplings of respective transition to the cavity
mode, ${\Omega}_{a}$ and ${\Omega}_{b}$ are the Rabi frequency of laser with
frequency $\omega_{a}$ and $\omega_{b}$.
Here we discuss very briefly about an effective $z$-component of interactions
(${{\sigma}_{i}}^{z}{{\sigma}_{i+1}}^{z}$) in such a system. The authors of
Ref.hart1 ; hart2 have proposed the same atomic level configuration but
having only one laser of frequency ${\omega}$ that mediates the atom-atom
coupling through virtual photons. Another laser field with frequency $\nu$ is
used to tune the effective magnetic field. In this case the Hamiltonian
${H_{AC}}$ changes but the Hamiltonians $H_{A}$ and $H_{C}$ are the same.
$\displaystyle{H_{AC}}$ $\displaystyle=$
$\displaystyle\sum_{j=1}^{N}[(\frac{{\Omega}}{2}e^{-i{{\omega}}t}+\frac{{\Lambda}}{2}e^{-i{{\nu}_{a}}t}{g_{a}}{a_{j}})|e_{j}><a_{j}|+h.c]$
(11) $\displaystyle+[a\leftrightarrow b].$
Here, ${\Omega}_{a}$ and ${\Omega}_{b}$ are the Rabi frequencies of the
driving laser with frequency ${\omega}$ on transition $|a>\rightarrow|e>$ ,
$|b>\rightarrow|e>$, whereas ${\Lambda}_{a}$ and ${\Lambda}_{b}$ are the
driving laser with frequency ${\nu}$ on transition $|a>\rightarrow|e>$ ,
$|b>\rightarrow|e>$. One can eliminate adiabatically the excited atomic levels
and photons by considering the interaction picture with respect to
$H_{0}=H_{A}~{}+~{}H_{C}$ [6,7]. They have considered the detuning parameter
in such a way that the Raman transitions between two level are suppressed and
also chosen the parameter in such a way that the dominant two-photon processes
are those that involve one laser photon and one cavity photon but the atom
makes no transition between levels a and b. Whenever two atoms exchange a
virtual photon both of them experience a Stark shift and play the role of an
effective ${{\sigma}^{z}}{{\sigma}^{z}}$ interaction hart1 ; hart2 ; suj2 .
Then the effective Hamiltonian reduces to
${H_{zz}}~{}=~{}\sum_{j=1}^{N}({B_{z}}{{\sigma}_{j}}^{z}~{}+~{}{J_{z}}{{\sigma}_{j}}^{z}{{\sigma}_{j+1}}^{z})$
(12)
These two parameters can be tuned independently by varying the laser
frequencies. Finally, they have obtained an effective model by combining
Hamiltonians $H_{xy}$ and $H_{zz}$ by using Suzuki-Trotter formalism. The
effective Hamiltonian simulated by this procedure is
$H_{spin}~{}=~{}\sum_{j=1}^{N}(B_{tot}{{\sigma}_{j}}^{z}~{}+~{}\sum_{{\alpha}=x,y,z}J_{\alpha}{{\sigma}_{j}}^{\alpha}{{\sigma}_{j+1}}^{\alpha})$
(13)
where $B_{tot}=B+{B_{z}}$. It has been shown in Ref. hart2 that $J_{y}$ is
less than $J_{x}$. From the analytical expressions of $J_{x}$ and $J_{y}$, it
is clear that the magnitudes of ${J_{1}}$ and $J_{2}$ are different. The
result of numerical simulations trigger us also to define a model to study the
quantum phases of this system. In the next section, we present the RG study of
this model Hamiltonian to extract quantum phases and transitions between them.
Here we present the analytical expression of $B_{tot}$, $J_{z}$ in terms of
different physical parameters of the system. We follow the references hart1 ;
james ; suj2 .
${J_{z}}={\gamma}_{2}{|\frac{{{\Omega}_{b}}^{*}g_{b}}{4\Delta_{b}}-\frac{{{\Omega}_{a}}^{*}g_{a}}{4\Delta_{a}}|}^{2}$
$B_{tot}=-\frac{1}{2}[\frac{{|{\Lambda}_{b}|}^{2}}{16{\tilde{\Delta_{b}}}^{2}}(4\tilde{\Delta_{b}}-\frac{{|{\Lambda_{a}}|}^{2}}{4({\tilde{\Delta}_{a}}-{\tilde{\Delta}_{b}})}-\frac{{|{\Lambda_{b}}|}^{2}}{{\tilde{\Delta}_{b}}}-{\beta_{2}})-{\beta_{3}}].$
${\beta_{2}}=\sum_{j=a,b}\frac{{|{\Omega_{j}}|}^{2}}{4({\Delta}_{j}-\tilde{{\Delta}_{b}})}.4\tilde{\gamma_{jb}}{g_{j}}^{2}$
$\displaystyle\beta_{3}$ $\displaystyle=$
$\displaystyle[\frac{{|{\Omega_{b}}|}^{2}}{16{{\Delta}_{b}}^{2}}(4{\Delta}_{b}-\frac{{|{\Omega_{a}}|}^{2}}{4{\Delta}_{b}}-$
(14)
$\displaystyle\frac{{|{\Omega_{b}}|}^{2}}{4({\Delta}_{a}-{\Delta}_{b})}-\frac{{|{\Omega_{b}}|}^{2}}{\Delta_{b}}-\sum_{j=a,b}\frac{{|{\Lambda_{j}}|}^{2}}{4({\Delta}_{j}-\tilde{{\Delta}_{b}})}.4{\gamma_{jb}}{g_{j}}^{2})+{\gamma_{bb}}^{2}\frac{{g_{b}}^{4}}{{\Delta_{b}}}-(a\leftrightarrow
b)]$
Here ${\gamma}_{1}=\frac{1}{N}\sum_{k}\frac{1}{\omega-{\omega}_{k}}$,
${\gamma}_{2}=\frac{1}{N}\sum_{k}\frac{e^{ik}}{\omega-{\omega}_{k}}$,
${\gamma}_{aa}={\gamma}_{bb}=\frac{1}{N}\sum_{k}\frac{1}{\omega-\omega_{k}}$.
${\gamma}_{ab}={\gamma}_{ba}=\frac{1}{N}\sum_{k}\frac{1}{\omega\pm\omega_{ab}-{\omega}_{k}}$
$\tilde{{\gamma}}_{ab}=\tilde{{\gamma}}_{ba}=\frac{1}{N}\sum_{k}\frac{1}{\nu\pm\omega_{ab}-{\omega}_{k}}$
$\tilde{{\gamma}}_{aa}=\tilde{{\gamma}}_{bb}=\frac{1}{N}\sum_{k}\frac{1}{\nu-\omega_{k}}$.
Analytical Derivation and Analysis of RG Equations:
To study the different quantum phases of the system described by the
Hamiltonian (Eq. 13), we express this Hamiltonian in more explicit way,
$\displaystyle H_{2}~{}=~{}\sum_{n}~{}[$
$\displaystyle(1+a)~{}S_{n}^{x}S_{n+1}^{x}~{}+~{}(1-a)~{}S_{n}^{y}S_{n+1}^{y}$
(15)
$\displaystyle+~{}\Delta~{}S_{n}^{z}S_{n+1}^{z}~{}+~{}h~{}S_{n}^{z}~{}]~{},$
where $S_{n}^{\alpha}$ are the spin-1/2 operators. We assume that the $XY$
anisotropy $a$ and the $zz$ coupling $\Delta$ satisfy the condition
$-1\leq\Delta\leq 1$, and $0<a\leq 1$ and magnetic field strength is $h\geq
0$. The parameters correspondence between the micro cavities and spin chain
are the following, $h\sim B_{tot}$, ${\Delta=J_{z}}$, ${J_{1}=1}$ and
${J_{2}}=a$. The $XY$ anisotropy breaks the in plane rotational symmetry. The
study of the quantum phases from the perspective of quantum spin system and
magnetism is not entirely a new one suj ; ric ; zamo . Here our main aim is to
study the quantum phases of microcavities array through the RG analysis of
this model Hamiltonian.
Spin operators can be recast in terms of spinless fermions through Jordan-
Wigner transformation and then finally one can express the spinless fermions
in terms of bosonic fields gia . We recast the spinless fermions operators in
terms of field operators by this relation.
${\psi}(x)~{}=~{}~{}[e^{ik_{F}x}~{}{\psi}_{R}(x)~{}+~{}e^{-ik_{F}x}~{}{\psi}_{L}(x)]$
, where ${\psi}_{R}(x)$ and ${\psi}_{L}(x)$ describe the second-quantized
fields of right- and the left-moving fermions respectively, and $k_{F}$ is the
Fermi wave vector. We express the fermionic fields in terms of bosonic field
by the relation
${{\psi}_{r}}(x)~{}=~{}~{}\frac{U_{r}}{\sqrt{2\pi\alpha}}~{}~{}e^{-i~{}(r\phi(x)~{}-~{}\theta(x))},$
where $r$ denotes the chirality of the fermionic fields, right (1) or left
movers (-1). The operators $U_{r}$ is the Klein factor to preserve the anti-
commutivity of fermions. $\phi$ field corresponds to the quantum fluctuations
(bosonic) of spin and $\theta$ is the dual field of $\phi$. They are related
by the relations ${\phi}_{R}~{}=~{}~{}\theta~{}-~{}\phi$ and
${\phi}_{L}~{}=~{}~{}\theta~{}+~{}\phi$. Hamiltonian
$H_{0}=\frac{v}{2}~{}\int~{}dx~{}[~{}(\partial_{x}\theta)^{2}~{}+~{}(\partial_{x}\phi)^{2}~{}]$
is non-interacting part of $H_{XYZ}$. Here $v$ is the velocity of the low-
energy excitations. It is one of the Luttinger liquid parameters and the other
is $K$, which is related to $\Delta$ by gia ; ric
$K~{}=~{}\frac{\pi}{\pi+2\sin^{-1}(\Delta)}~{}.$ (16)
where $K$ takes the values 1 and 1/2 for $\Delta=0$ (free field), and
$\Delta=1$ (isotropic anti-ferromagnet), respectively. The relation between
$K$ and $\Delta$ is not preserved under the renormalization, so this relation
is only correct for the initial Hamiltonian. The analytical form of the spin
operators in terms of the bosonic fields are:
$S_{n}^{x}~{}=~{}[~{}c_{2}\cos(2{\sqrt{\pi
K}}\phi)~{}+~{}(-1)^{n}c_{3}~{}]~{}\cos({\sqrt{\frac{\pi}{K}}}\theta)$;
$S_{n}^{y}~{}=~{}-[~{}c_{2}\cos(2{\sqrt{\pi
K}}\phi)~{}+~{}(-1)^{n}c_{3}~{}]~{}\sin({\sqrt{\frac{\pi}{K}}}\theta)$, and
$S_{n}^{z}~{}=~{}{\sqrt{\frac{\pi}{K}}}~{}\partial_{x}\phi~{}+~{}(-1)^{n}c_{1}\cos(2{\sqrt{\pi
K}}\phi)~{}$ where $c_{i}$’s are constants as given in Ref. zamo . The
Hamiltonian $H_{2}$ in terms of bosonic fields is the following,
$\displaystyle H_{2}$ $\displaystyle=$ $\displaystyle
H_{0}+\frac{a}{2\pi\alpha}\int\cos(2{\sqrt{\frac{\pi}{K}}}\theta(x))dx$ (17)
$\displaystyle+\frac{\Delta}{{2\pi\alpha}^{2}}\int\cos(4{\sqrt{\pi
K}}\phi(x))dx+\frac{h\sqrt{K}}{\pi\alpha}\int{{\partial}_{x}}{\phi(x)}dx$
One can also write the above Hamiltonian in the following form
$\displaystyle H_{2}$ $\displaystyle=$ $\displaystyle
H_{01}+\frac{a}{2\pi\alpha}\int\cos(2{\sqrt{{\pi}}}\theta(x))dx$ (18)
$\displaystyle+\frac{\Delta}{{2\pi\alpha}^{2}}\int\cos(4{\sqrt{\pi}}\phi(x))dx+\frac{h}{\pi\alpha}\int{{\partial}_{x}}{\phi(x)}dx$
Where $H_{01}$,
$\displaystyle H_{01}=\frac{1}{2\pi}\int
dx[uK{({\nabla\theta(x)})}^{2}+(u/K){({\nabla\phi(x)})}^{2}]$ (19)
One can get the $H_{XY}$ Hamiltonian by simply putting $\Delta=0$ in the above
Hamiltonian. In this derivation, different powers of coefficients $c_{i}$ have
been absorbed in the definition of $a,h$ and $\Delta$. The integration of the
oscillatory terms in the Hamiltonian yield negligible small contributions and
the origin of the oscillatory terms occur due the spin operators. So it’s a
reasonably good approximation to keep only the non-oscillatory terms in the
Hamiltonian. The Gaussian scaling dimension of these coupling terms, $a$ and
$\Delta$ are $1/K$ and $4K$ respectively. The third term ($\Delta$) of the
Hamiltonian tends to order the system into density wave phase , whereas the
second term ($a$) of the Hamiltonian favors the staggered order in the $XY$
plane. Two sine-Gordon coupling terms are from two dual fields. Therefore, the
model Hamiltonian consists of two competing interactions. This Hamiltonian
contains two strongly relevant and mutually nonlocal perturbation over the
Gaussian (critical) theory. In such a situation the strong coupling fixed
point is usually determined by the most relevant perturbation whose amplitude
grows up according to its Gaussian scaling dimensions and it is not much
affected by the less relevant coupling terms. However, this is not the general
rule if the two operators exclude each other, i.e., if the field
configurations which minimize one perturbation term do not minimize the other.
In this case interplay between the two competing relevant operators can
produce a novel quantum phase transition through a critical point or a
critical line. Therefore, we would like to study the RG equation to interpret
the quantum phases of the system. In the RG theory, we not only able to
predict the weak coupling limit but also the the intermediate values of the
coupling. The RG theory is a perturbative theory and it ceases to be valid
when the coupling constant $g(l)\sim 1$.
We now study how the parameters $a$, $\Delta$ and $K$ flow under RG. The
operators in Eq. (17) are related to each other through the operator product
expansion. So the RG equations for their coefficients therefore are coupled to
each other. We use operator product expansion to derive these RG equations
which is independent of boundary condition cardy . In our derivation, we
consider two operators, $X_{1}=e^{(ia_{1}\phi+ib_{1}\theta)}$ and
$X_{2}=e^{(ia_{2}\phi+ib_{2}\theta)}$. In the RG procedure, one can write
these two field operators as a sum of fast and slow mode fields. In the fast
field, the momentum range is $\Lambda e^{-dl}<K<\Lambda$ and for the slow
field $K<\Lambda e^{-dl}$, where $\Lambda$ is the momentum cut-off, $dl$ is
the change in the logarithmic scale. The next step is the integration of the
fast field for the operators $X_{1}$ and $X_{2}$, it yields a third operator
at the same space time point,
$X_{3}=e^{i(a_{1}+a_{2})\phi+i(b_{1}+b_{2})\theta)}$. The prefactor of $X_{3}$
can be found by the relation, ${X_{1}}{X_{2}}\sim
e^{-(a_{1}a_{2}+b_{1}b_{2})}\frac{dl}{2\pi}X_{3}$. Our Hamiltonian consists of
two operators, if we consider $l_{1}$ and $l_{2}$ as the coefficient of the
operators $X_{1}$ and $X_{2}$ respectively. Then the RG expressions for
$\frac{dX_{3}}{dl}$ contains the term
$(a_{1}a_{2}+b_{1}b_{2})\frac{l_{1}l_{2}}{2\pi}$. This is the procedure to
derive these RG equations.
In the RG process, one can write RG equations themselves are established in a
perturbative expansion in coupling constant ($g(l)$). They cease to be valid
beyond a certain length scale, where $g(l)\sim 1$ gia . The RG equations for
the coefficients of Hamiltonian $H_{XYZ}$ are
$\displaystyle\frac{da}{dl}~{}$ $\displaystyle=$
$\displaystyle~{}(2-\frac{1}{K})a,$ $\displaystyle\frac{d{\Delta}}{dl}~{}$
$\displaystyle=$ $\displaystyle~{}(2-4K){\Delta}$
$\displaystyle\frac{dK}{dl}~{}$ $\displaystyle=$
$\displaystyle~{}\frac{a^{2}}{4}~{}-~{}K^{2}{\Delta}^{2}~{},$ (20)
We have followed Ref.suj during the derivation of these RG equations. These
RG equations have trivial (${a^{*}}=0={{\Delta}^{*}}$) fixed points for any
arbitrary $K$. Apart from that these RG equations have also two non-trivial
fixed lines, $a=\Delta$ and $a=-{{\Delta}}$ for $K=1/2$. The above RG
equations show that there is no duality in flow diagram. Here we mean,
duality, that if $\theta$ and $\phi$ interchange $\theta\leftrightarrow\phi$,
$K\leftrightarrow K^{-1}$ and $\Delta\leftrightarrow a$ will not produce the
same set of RG equations.
For $K>1/2$, the sine-Gordon coupling term correspond to the anisotropic
exchange coupling become relevant and the system flows to the photonic
superfluid phase. For $K<1/2$, the sine-Gordon coupling term correspond to the
z-component exchange interaction become relevant and the system flowing off to
the photon blocked induced insulating phase.
Here we explain the physical significance of different quantum phases of the
atom-cavity system what we find in our study.
(1). When both the anisotropic exchange interaction and the z-component of
exchange interaction is absent then there are no sine-Gordon coupling terms in
the Hamiltonian. Then the system is in the mass less Luttinger liquid phase,
i.e., the system is in the photonic metallic state. The other source of
photonic metallic state is that when the system shows the BKT transition which
we will discuss in the next section.
(2). The photonic insulator state of the atom-cavity array system corresponds
to the insulating state of the system where there is no transmission of photon
between the microcavities in the array due to the interaction between photons.
The appearance of this phase has already discussed in the analysis of three RG
equations. This phase will occurs when the system shows the BKT transition,
which we will discuss in the next section.
(3). Photonic superfluid state of the atom-cavity array system corresponds to
the gapless excitations of the system where the photon transmit from one
cavity to the other without any blocking . For this one dimensional cavity QED
system where there is no order parameter, one can only discriminate between
the photonic metallic state and photonic superfluid state by only finding the
difference of fluctuation in photon number in every sites of the array. This
photonic superfluid phase corresponds to dissipitionless flow of photon in the
system. We have already discussed about the appearance of photonic superfluid
phase from the analysis of three RG equations. The other source of the
appearance of photonic superfluid phase is the BKT transition which we will
discuss in the next section.
Berezinskii-Kosterlitz-Transition Physics in Cavity QED System.
The physics of BKT transition has found in different one dimensional and two
dimensional (classical system) and it has discussed extensively in different
context in the Ref. subir ; gia ; druf .
Before we start to discuss the appearance of BKT transition in our system, we
would like to discuss very briefly why it is necessary to study the BKT
transition. Here we study two different situations of our model Hamiltonian.
For the first case the exchange anisotropy is absent ($J_{2}=0$) and for the
second case z-component of exchange interaction is absent ($J_{z}=0$). For
both of these cases only one of the sine-Gordon coupling term is present,
therefore, there is no competition between the two mutually non local
perturbation. Therefore one can think that there is no need to study the RG to
extract the quantum phases and phase boundaries. But we still apply RG method
for the following reason. Each of these Hamiltonians consist of two part, the
first one ($H_{01}$ ) is the non-interacting where the $\phi$ and $\theta$
fields show the quadratic fluctuations and the other part of these
Hamiltonians are the sine-Gordon coupling terms which of either $\theta$ or
$\phi$ fields. The sine-Gordon coupling term lock the field either $\theta$ or
$\phi$ in the minima of the potential well. Therefore the system has a
competition between the quadratic part of the Hamiltonian and the sine-Gordon
coupling term and this competition will govern the low energy physics of these
Hamiltonians in different limit of the system. The RG process (BKT transition)
will predict the quantum phases of these system correctly.
It is very clear from the analytical expression of ${J_{1}},J_{2}$ and $J_{z}$
that one can control these parameters in the laboratory. The quantum state of
engineering of cavity QED system is in the state of art due to the rapid
technological development of this system [1]. In this study we consider the
situations where the $J_{2}$ is absent, i.e., the system with a single Rabi
frequency oscillation. We also consider the situation where the z-component of
exchange interaction is absent. These two situations lead to the two different
set of RG equations which show BKT transition.
For the first case, there is no anisotropic exchange coupling. In this
situation, the three RG equations of the previous section reduce to two RG
equations, which are the following
$\displaystyle\frac{d{\Delta}}{dl}~{}$ $\displaystyle=$
$\displaystyle~{}(2-4K){\Delta}$ $\displaystyle\frac{dK}{dl}~{}$
$\displaystyle=$ $\displaystyle~{}-~{}K^{2}{\Delta}^{2}~{},$ (21)
Now we express these RG equations in the form of BKT transition form. As we
understand from the RG equations that the transition occurs at $K=1/2$. To
study the flow of the RG equation around this transition point, we recast the
RG equation in suitable form.
Here we follow the following transformation, $K=1/2+\frac{y_{||}}{4}$ and
$\Delta\rightarrow\Delta/2$.
The above equation reduce to the standard BKT equation.
$\displaystyle\frac{dy_{||}}{dl}~{}$ $\displaystyle=$
$\displaystyle~{}-{\Delta}^{2}$ $\displaystyle\frac{d\Delta}{dl}~{}$
$\displaystyle=$ $\displaystyle~{}-~{}y_{||}{\Delta}~{},$ (22)
In our case, $\Delta\frac{d\Delta}{dl}=y_{||}\frac{dy_{||}}{dl}$. Therefore
$X^{2}={y_{||}}^{2}-{\Delta}^{2}$ is a constant of motion. Here we discuss the
relevant physics and the quantum phase transition between the photonic
metallic state and photonic Coulomb blocked state based on these equations.
Here we consider the following situations based on these equations:
(1) When $X>0$ and $y_{||}>0$, the sine-Gordon coupling term corresponding to
$\Delta$ term is irrelevant, the fixed point, ${\Delta}^{*}=0$ and
${y_{||}}^{*}=X$, close to this fixed point, we can write the RG equations as
$\frac{d\Delta}{dl}=2(1-2K^{*}){\Delta(l)}$ and $\frac{dy_{||}}{dl}=0$. Using
the flow equation and constant of motion for ${y_{||}}>{\Delta}$. One can
write the solution of coupling terms as
$y_{||}(l)=\frac{X}{tanh(Xl+atanh(\frac{X}{{y_{||}}^{0}})}$ (23)
$\Delta(l)=\frac{X}{sin(Xl+atanh(\frac{X}{{y_{||}}^{0}})}$ (24)
The condition at the line of seperatix is
${y_{||}}(l)={\Delta(l)}=\frac{{\Delta}^{0}}{1+{\Delta}^{0}l}$. There is no
mass gap excitation in the elementary excitation of the system, i.e, there is
no photon blocked induced insulating phase. Hence the system is in the
photonic metallic state.
To the best of our knowledge for the first time in the literature we predict
the existence of photonic metallic state for the array of Cavity QED system.
Here we would like to present the basic origin of photonic-metallic state
explicitly.
If we do the Jordan-Wigner transformation of our model Hamiltonian as we
present in Ref. 22. It is very clear from the Eq. 32 of Ref. 22 that the first
term represent the photon hopping term across the lattice of the cavity QED
array. The second and third term of the Hamiltonian represent the photonic
pair correlation and photonic density wave respectively in the system. When
the second and third term in the Hamiltonian either are absent or irrelevant
in the RG sense at that situation the system shows the photonic metallic
state.
(2) When ${\Delta}>y_{||}$. The RG equation for the coupling $\Delta(l)$
flowing off to the strong coupling regime. A perturbative expansion in
$\Delta$ cease to be valid beyond a certain length scale for which
$\Delta(l)\sim 1$. The analysis of the RG equation is not valid beyond this
length scale.
The analytical relation between the coupling constant is
$arctan(\frac{{{y}_{||}}^{0}}{\sqrt{{{\Delta}_{0}}^{2}-{{y_{||}}^{0}}^{2}}})-arctan(\frac{{{y}_{||}}}{\sqrt{{{\Delta}_{0}}^{2}-{{y_{||}}^{0}}^{2}}})=\sqrt{{{\Delta}_{0}}^{2}-{{y_{||}}^{0}}^{2}}.$
(25)
In this limit, the RG flowing off to the strong coupling phase. Here we
discuss the relevant physics of the phase. The sine-Gordon coupling term is
$\frac{\Delta_{1}u}{2\pi{\alpha}^{2}}\int dxcos(4\phi(x)).$
where $\Delta_{1}=\frac{\Delta}{\pi u}$. The $\phi(x)$ field locks into one of
the minima of the cosine potential. Now we expand the potential for large
$\Delta_{1}$ in the spirit of usual RG method subir ; gia ; druf . We can
write the effective Hamiltonian near to the minima
$H=H_{0}+\frac{4\Delta_{1}u}{\pi{\alpha}^{2}}{\phi}^{2}(x)$
. The total action of the Hamiltonian of the system can be written as
$S=\frac{1}{2\pi
K}\frac{1}{\beta\Omega}\sum_{k,{\omega}_{n}}[\frac{{{\omega}_{n}}^{2}}{u}+uk^{2}+\frac{8K\Delta_{1}u}{{\alpha}^{2}}]{\phi}^{*}(k,\omega_{n}){\phi}(k,\omega_{n})$
(26)
The excitations of the system which cost a finite energy even at $k=0$. In
this situation the field, $\phi(x)$ is massive. The system posses phononic
type mode. This phononic mode is the small oscillation of the field
${\phi}(x)$ around the minima of the cosine potential.
Suppose we consider the RG equation up to the point where ${\Delta_{1}}(l)\sim
1$. This excitations gap in the spectrum has the dimension of an energy and
the renormalization relation is ${\Omega}_{M}(l)=e^{l}{\Omega}_{M}(l=0)$, for
the case, ${\Delta_{1}}(l^{*})\sim 1$. We can write the expression for the gap
following the action as,
${\Omega}_{M}(l^{*})\sim\sqrt{{\Delta_{1}(l^{*})}}u/{\alpha}$. The true gap of
the system is ${\Omega}_{M}(l=0)\sim e^{-l^{*}}{\Omega}_{M}(l^{*})$,
${\Omega}_{M}(l^{*})=\frac{{u}}{\alpha}$. Thus at this phase of this system,
the explicit dependence of the Luttinger liquid parameter is absent. This gap
of the system is of the order of bandwidth of the cavity QED system.
(3) Now we consider the case, when ${\Delta_{1}}<<|y_{||}|$. The system is in
the deep massive phase. In this phase, one can write the RG equation as
$\frac{dK(l)}{dl}=0$ and $\frac{d\Delta_{1}(l)}{dl}=2(1-2K){\Delta_{1}}(l)$.
For this RG equation, we get ${\Delta_{1}}(l)={\Delta_{1}}(0)e^{2(1-2K)l}$
from this equation $e^{-l^{*}}={\Delta(0)}^{1/2(1-2K)}$. The true gap of the
system is ${\Omega}_{2}(l=0)\sim e^{-l^{*}}{\Omega}_{2}(l^{*})$,
$\frac{\Omega_{2}(l=0)}{\Omega_{2}(l^{*})}\sim{\Delta_{1}(0)}^{+1/2(1-2K)}$.
This gap is the power law dependence of the bare $\Delta_{1}$. This gap varies
with $K$, i.e., the gap of the system is now varying with the interaction of
the system. It is very clear from the analytical expression that the gap gets
smaller and tends to zero at $K=1/2$, this prediction is consistent with the
physical scenario that at $K=1/2$ the system shows the phase transition.
(4) Close to the transition point, Eq.(25) reduces to
$\sqrt{{{\Delta_{1}}_{0}}^{2}-{(y_{||})}^{2}}{l^{*}}=\pi$. So the square root
term goes to zero at the transition point and therefore
${y_{||}}^{0}/{\sqrt{{{\Delta_{1}}_{0}}^{2}-{y_{||}}^{2}}}\rightarrow\infty$
and the gap
${\Omega}(l=0)\sim{\Omega}({l^{*}})e^{-\frac{\pi}{\sqrt{{{\Delta_{1}}_{0}}^{2}-{y_{||}}^{2}}}}$.
As one approach to the transition point, such as
${\Delta_{1}}_{0}\rightarrow|{y_{||}}|$. The gap is exponentially small in the
square root of the distance to the transition.
It is therefore clear from our RG analysis that the system is in the massive
phase for the two limits. For the first case when ${\Delta}>>{y_{||}}$, the
excitation gap of the system is of the order of the bandwidth of the system
and the effect of photonic strong correlation is not explicit. The other limit
of gapped state, i.e, when $\Delta<<|y_{||}|$. The excitation gap of the
system has a power law dependence which vary with $K$.
It is very clear from the analytical expression of $J_{z}$ that one can tune
it to zero by adjusting the Rabi frequencies oscillation of the system. In
this situation, there is no z-component of intercavity exchange interaction.
Therefore the three RG equations are reduce to two RG equations.
$\displaystyle\frac{da}{dl}~{}$ $\displaystyle=$ $\displaystyle~{}(2-1/K)a$
$\displaystyle\frac{dlnK}{dl}~{}$ $\displaystyle=$
$\displaystyle~{}~{}{a^{2}}/{4K}~{},$ (27)
If we do the following transformation $\tilde{K}=1/{2K}$ and $\tilde{a}=a/2$
and after that we follow another transformations $\tilde{K}=1+{y_{||}}/2$ the
above equations reduce to
$\displaystyle\frac{dy_{||}}{dl}~{}$ $\displaystyle=$
$\displaystyle~{}-{\tilde{a}}^{2}$ $\displaystyle\frac{d\tilde{a}}{dl}~{}$
$\displaystyle=$ $\displaystyle~{}-~{}y_{||}{\tilde{a}}~{},$ (28)
The mathematical structure of these equations are the same as that of Eq. 22.
Therefore the mathematical analysis for the different limits are the same but
massive phase and deep massive phase for these situations are the photonic
superfluid phase.
The sine-Gordon coupling term is
$\frac{a_{2}u}{2\pi{\alpha}}\int dxcos(2\theta(x)).$
where $a_{2}=\frac{a}{u}$. The $\theta(x)$ field locks into one of the minima
of the cosine potential. Now we expand the potential for large $a_{2}$. We can
write the effective Hamiltonian near to the minima
$H=H_{0}+\frac{a_{2}u}{{\pi\alpha}}{\theta}^{2}(x).$
The total action of the Hamiltonian of the system can be written as
$S=\frac{1}{2\pi
K}\frac{1}{\beta\Omega}\sum_{k,{\omega}_{n}}[\frac{{{\omega}_{n}}^{2}}{u}+uk^{2}+\frac{2Ka_{2}u}{{\alpha}}]{\theta}^{*}(k,\omega_{n}){\theta}(k,\omega_{n})$
(29)
The excitations of the system which cost a finite energy even at $k=0$. In
this situation the field, $\theta(x)$ is massive. The true gap of the system
is ${\Omega}_{C}(l=0)\sim e^{-l^{*}}{\Omega}_{C}(l^{*})$,
${\Omega}_{C}(l^{*})=(\frac{{u}}{\alpha})\sqrt{\alpha}$. Thus at this phase of
this system, the explicit dependence of the Luttinger liquid parameter is
absent. This gap of the system is much less than the excitation gap of the
photon blocked state due to the presence of the extra factor $\sqrt{\alpha}$.
In the deep massive phase,
$\frac{\Omega_{C}(l=0)}{\Omega_{C}(l^{*})}\sim{a_{2}(0)}^{+1/(2-1/K)}$. This
gap is the power law dependence of the bare $a_{2}$. This gap varies with $K$,
i.e., the gap of the system is now varying with the interaction of the system.
It is very clear from the analytical expression that the gap gets smaller and
tends to zero at $K=1/2$ and this prediction is consistent with the physical
scenario that at $K=1/2$ the system shows the phase transition from photonic
metallic state to photonic superfluid phase.
Now we consider the effect of magnetic field ( ${\delta}_{1}$ term) on the
quantum phases and phase boundaries of this cavity QED arrays.
We do the following transformation to eliminate the magnetic field term from
the Hamiltonian. We substitute $4\pi\sqrt{K}\phi\rightarrow
4\pi\sqrt{K}\phi+4\pi\delta x$, where $\delta=\sqrt{K}h$. This introduce the
spatial oscillation in the $\Delta$ term. When the coupling $a$ is absent and
the applied detuning field is larger than the photonic Mott gap, the system
drives to the gapless photonic metallic state. This is nothing but the well
known Porkovosky-Talapov model which shows the commensurate to incommensurate
transition gia ; suj . When $\delta{a_{1}}>>1$ (where $a_{1}$ is the lattice
spacing), the $cos(4\pi\sqrt{K}(\phi+\delta x))$ term is quickly oscillating
and averages out zero. Thus the system reflects the competition between the
$\Delta$ and $h$. As a results of it, the RG flow for $\Delta$ has to be
cutoff when $4\pi\sqrt{K}\delta(l)a_{1}\sim 1$. To the lowest order in
$\Delta,a$ and $\delta$, the RG flow equation is
$\frac{d\delta}{dl}=\delta$ (30)
As a result of this RG equation, the system shows two different response,
either the flow ${\Delta}_{1}$ and $\delta$ or $a$ and $\delta$ flows of to
the strong coupling phase.
When $\Delta_{1}(l)$ reaches strong coupling before $\delta a_{1}$ become
order one. The phase boundary is the same as we predict without magnetic
field. The condition when $\Delta_{1}(l^{*})\sim 1$ define a new length scale
$l^{*}$. It generates a self-consistent scenario that
$4\pi\delta(l^{*})a_{1}<<{{\Delta_{1}}(0)}^{1/(2-4K)}$. In the other limit
when the RG flow of $a$ flows off to the strong coupling phase does not
affected by the magnetic field term and thus the system is in the photonic
superfluid phase.
Conclusions: We have presented three sets of RG equations for the different
physical situations for the Cavity QED system. We have predicted two different
BKT transitions for the different physical situations. We have predicted the
photonic superfluid phase, photonic metallic phase and photon blocked induced
insulating phase. To the best of our knowledge, this explicit quantum phase
analysis and their behavior based on the RG study for the microcavity array
system is absent in the previous literature.
Acknowledgments
The author would like to acknowledge the series of lectures and also
discussions of Prof. S. M. Girvin during the International Workshop/School on
Dirac Materials and Chandrashekar Discussion Meeting, December’2012. The
author would like to thank The Center for Condensed Matter Theory of the
Physics Department of IISc for extended facility. Finally the author would
like to thank Prof. Prabir Mukherjee for reading this manuscript critically.
## References
* (1) G. S. Agarwal, Quantum Optics (Cambridge University Press, Cambridge 2013).
* (2) S. M. Girvin, Circuit QED: Superconducting Qubits Coupled to Microwave Photons (Lecture Notes/Private Communication).
* (3) Hartmann Michael J, Fernando G S, Brando L and Plenio Martin B 2006 Nature Phys 462 849; Hartmann Michael J, Fernando G S, Brando L and Plenio Martin B 2008, Laser and Photonics Rev. 2 527\.
* (4) Hartmann Michael J, Fernando G S, Brando L and Plenio Martin B 2007, Phys. Rev. Lett 99 160501\.
* (5) S. Sarkar, Physica B 407, 44 (2012).
* (6) S. Horoche and J. M. Raimond in Exploring the Quantum Atoms, Cavities, and Photons, (Oxford University Press, 2006).
* (7) E. T. Jaynes and F. W. Cummings, Proc IEEE 51, 89 (1963).
* (8) A. C. Ji, X. C. Xie, and W. M. Liu, Phys. Rev.Lett. 99, 183602 (2007).
* (9) T. Byrnes, N. Y. Kim, K. Kusudo, and Y. Yamamoto, Phys. Rev. B 78 075320 (2008).
* (10) I. Carusotto $et~{}al.$, arXiv:0812.4195.
* (11) M. J. Bhaseen $et~{}al.$, Phys. Rev. Lett. 102 135301 (2009).
* (12) J. Zhao, A. W. Sandvik and K. Ueda, arXiv:0806.3603.
* (13) P. Pippan, H. G. Evertz, and M. Hohenadler, arXiv: 0904.1350.
* (14) M. Aichhorn $et~{}al.$, Phys. Rev. Lett. 100 216401 (2008).
* (15) S. Schmidt, and G. Blatter, arXiv:0905.3344 .
* (16) D. G. Angelakis, M. F. Santos and S. Bose, Phys. Rev. A 76 R031805, (2007).
* (17) B. Kumar, J. Phys. A: Math. Theor 42, 245307 (2009);
* (18) Sachdev Subir in ”Quantum Phase Transition” (Cambridge University Press, Cambridge, 1998).
* (19) A. K. Pal and I. Bose, J. Phys. B: At. Mol. Opt 44, 045101 (2011).
* (20) T. Giamarchi in Quantum Physics in One Dimension (Clarendon Press, Oxford, 2004).
* (21) Itzykson and Drouffe in Statistical Field Theory, Volume 1 (Cambridge University Press, New York 1989).
* (22) Here we present the Jordon-Wigner transformation of Eq. 15, $Q_{i}=e^{i\pi\sum_{n=1}^{i-1}{c_{i}}^{\dagger}c_{i}}$, ${S_{i}}^{-}=Q_{i}c_{i}$, ${S_{i}}^{+}={c_{i}}^{\dagger}{Q_{i}}^{\dagger}$. The Eq. 15 under this trnsformation become
$\displaystyle H$ $\displaystyle=$
$\displaystyle\sum_{i=1}^{n}({c_{i+1}}^{\dagger}c_{i}+h.c)+a({c_{i}}^{\dagger}{c_{i+1}}^{\dagger}+h.c)$
(31)
$\displaystyle+\Delta\sum_{i}({c_{i}}^{\dagger}c_{i}-1/2)({c_{i+1}}^{\dagger}c_{i+1}-1/2)+h({c_{i}}^{\dagger}c_{i}-1/2)$
It is clear from the above expression that the anisotropic-exchange
interaction induce creation or annihilation of photon pairs.
* (23) K. M. Birnbaum $et~{}al.$, Nature 436, 87 (2005).
* (24) F. V. James and J. Jerke, Can. J. Phys 85, 625 (2007).
* (25) S. Sarkar, Phys. Rev. B 74, 052410 (2006).
* (26) U. Schollwock, J. Richter, D. J. J. Farnell, and R. F. Bishop (Eds.), Lect. Notes Phys. 645, ”Quantum Magnetism” (Springer, Berlin, 2004).
* (27) S. Lukyanov and A. Zamolodchikov, Nucl. Phys. B 493, 571 (1997).
* (28) J. Cardy in Scaling and Renormalization in Statistical Physics (Cambridge University Press, Cambridge, 1996); Affleck I 1989 in Fields, Strings and and Critical Phenomena, ed Brezin E and Zinn-Justin J (North-Holland, Amstardam).
|
arxiv-papers
| 2013-09-30T09:52:49 |
2024-09-04T02:49:51.706002
|
{
"license": "Public Domain",
"authors": "Sujit Sarkar",
"submitter": "Sujit Sarkar",
"url": "https://arxiv.org/abs/1309.7777"
}
|
1309.7912
|
# An Image-Based Fluid Surface Pattern Model
Mauro de Amorim\- [email protected] Ricardo Fabbri\- [email protected],br
Lucia Maria dos Santos Pinto\- [email protected] Francisco Duarte Moura
Neto\- [email protected] Polytechnic Institute at the Rio de Janeiro State
University 28630-050 - Nova Friburgo, RJ, Brazil
###### Abstract
This work aims at generating a model of the ocean surface and its motion from
one or more video cameras. The idea is to model wave patterns from video as a
first step towards a larger system of photogrammetric monitoring of marine
conditions for use in offshore oil drilling platforms. The first part of the
proposed approach consists in reducing the dimensionality of sensor data made
up of the many pixels of each frame of the input video streams. This enables
finding a concise number of most relevant parameters to model the temporal
dataset, yielding an efficient data-driven model of the evolution of the
observed surface. The second part proposes stochastic modeling to better
capture the patterns embedded in the data. One can then draw samples from the
final model, which are expected to simulate the behavior of previously
observed flow, in order to determine conditions that match new observations.
In this paper we focus on proposing and discussing the overall approach and on
comparing two different techniques for dimensionality reduction in the first
stage: principal component analysis and diffusion maps. Work is underway on
the second stage of constructing better stochastic models of fluid surface
motion as proposed here.
Keywords: Inverse problems, fluid motion, nonrigid 3D reconstruction,
dimensionality reduction, diffusion maps, pattern theory
## 1 Introduction
The simulation of fluids is an imporant tool in computer graphics, e.g., for
generating realistic animations of water flow. Many fluid simulations perform
the evolution of liquid through time based on the Navier-Stokes equation or
simpler wave and frequency models, among others. This work is an initial part
of a larger effort to perform the inverse problem of generating graphics
simulations, i.e., to automatically extract 3D models of fluid surfaces
starting from real-world video data. The obtained model can then be used for
simulating the motion of ocean surface patterns as _observed_ in the real
world. Moreover, this simulation can be matched to video at new time instants
in order to predict and track the actual conditions of the _observed_ fluid.
In order to track the apparently complex dynamics of a large number of images
of the ocean surface, our proposed system reduces the dimensionality of
previsouly observed data and automatically learns a concise data-driven model.
This enables the inference through synthesis of patterns that are intrisic to
the observed phenomenon, while drastically reducing the search space with
little to no loss. The remaining error in the low dimensional modeling is
dealt with as part of a stochastic modeling stage.
While detailed techniques for the plausible inference involving such patterns
is ongoing work, we propose that the patterns in newly observed data are to be
found by stochastic modeling, using the Bayesian paradigm, a method which can
be described as “analysis by synthesis” (Mumford & Desolneux,, 2010). Finally,
a stochastic model also enables generating plausible new images for a
realistic image-based simulation. Figure 1 shows a diagram representing the
steps of the proposed modeling approach dealt with in this manuscript.
Two dimensionality reduction techniques were assessed: _(i)_ principal
components analysis, which consists in performing orthogonal projections of
data along linear directions of greater variance; and _(ii)_ diffusion maps, a
recent non-linear technique which organizes data into graphs of local
similarity and extracts global structure through a diffusion process (Lafon &
Lee,, 2006).
Figure 1: Diagram of the proposed approach for the modeling of fluid surface
patterns aiming at applications of monitoring ocean conditions for offshore
oil drilling processes, among others.
## 2 Dimensionality Reduction Techniques
Assume a video is made up of $n$ video frames represented by a column vector
$y_{i}$, $i=1,\dots,n$. Each video frame is an image with a total of $p$
pixels, so that each observation or data point is a $p$-dimensional vector
$y_{i}\in\mathbbm{R}^{p}$. For instance, frames of a typical high-definition
video $1920\times 1080$ are each comprised of more than 2 milion dimensions,
an excessive number of parameters to capture the observable modes of variation
of fluid surface patterns that are relevant for monitoring applications. We
explore two different techniques for automatically reducing this number of
parameters as a step towards generating a tractable yet meaningful model.
### 2.1 Principal Component Analysis
The principal component analysis (PCA) is a linear dimentionality reduction
technique. High-dimensional data, i.e., $p$-dimensional, gets reduced by PCA
to a global linear submanifold of reduced dimension $q\ll p$, determined by
the directions of greatest variability in the data. These directions are given
by $q$ orthonormal vectors, the principal components, which are the
eigenvectors corresponding to the largest eigenvalues of the covariance matrix
of the data.
Principal component analysis can be computed through the singular value
decomposition (SVD) (Hastie & et al.,, 2013; Jolliffe,, 2002; Golub & Van
Loan,, 1996) as follows. The $n$ observations $y_{i}\in\mathbbm{R}^{p}$ are
mean-centered by subtracting their average
$\bar{y}=\dfrac{1}{n}\sum^{n}_{j=1}y_{j}$. The new data vectors
$x_{i}=y_{i}-\bar{y}$ form the columns of the $p\times n$ data matrix $X$
whose singluar value decomposition is given by
$X=U\Sigma V^{\top},$ (1)
with $U$ and $V$ being $p\times n$ and $n\times n$ orthonormal matrices,
resp., and $\Sigma$ an $n\times n$ diagonal matrix, with diagonal elements
$\sigma_{1}\geq\dots\geq\sigma_{n}\geq 0$ known as singular values. The first
$q$ columns of $U$ form the $p\times q$ matrix $U_{q}$, the so-called first
$q$ principal components, and $\alpha_{i}=U_{q}\top x_{i}$ are the coordinates
of any given observation $x_{i}$ in the subspace of the principal components.
The computation of PCA using SVD is equivalent to computing the eigenvalues
and eigenvectors of the covariance matrix, which is given by
$\dfrac{1}{n-1}XX^{\top}.$ (2)
It is possible to show this fact by taking the SVD of the matrix $X$ as shown
in (1), and building the product $XX^{\top}$, giving
$XX^{\top}=(U\Sigma V^{\top})(V\Sigma U^{\top}).$ (3)
Since $V$ is an orthogonal matrix, we have
$XX^{\top}=U\Sigma^{2}U^{\top}.$ (4)
From Equation (4) we clearly see the correspondence of the SVD and the
covariance of the data samples. The singular values of $X$ are the square
roots of the eigenvalues of the covariance matrix, and the singular vectors of
$X$ are the eigenvectors of the covariance matrix.
#### 2.1.1 Stochastic SVD
In our application, the dataset is very large and in high dimensions, in which
case computing the eigenvalues of the data matrix is a key challenge. Recent
techniques enable the computation of large matrix decompositions in a robust
manner without explicitly forming the entire data matrix in memory (Halko et
al.,, 2010, 2011). The key main idea behind these methods is stochasticity.
For the construction of the SVD using stochasticity, a method we refer in this
manuscript as stochastic SVD is available (Halko et al.,, 2011) which can be
broken in two stages. First, the construction of a subspace with reduced
dimension in which to represent the input data in an approximate manner; a
matrix is bult in which the columns form a reduced orthonormal basis for the
data. Second, the projection of the observations onto that subspace and the
subsequent computation of the SVD. Both stages are outlined in Algorithm 1.
Algorithm 1 – Stochastic SVD
Given an ${m\times n}$ matrix $A$, and integers $l$ and $q$, compute an
approximate decomposition of $A\approx U\Sigma V^{\top}$, where $U$ and $V$
are orthonormal and $\Sigma$ is a non-negative diagonal matrix.
1:Create an ${n\times l}$ matrix $\Omega$, with independent entries and
standard normal distribution.
2:Form $Y=A\Omega$.
3:Perform the $QR$ decomposition of $Y$, $Y=QR$, with $Q$ an $m\times l$
orthonormal matrix.
4:Form $B=Q^{\top}A$.
5:Compute an SVD of the reduced matrix $B=\tilde{U}\Sigma V^{\top}$.
6:Form the otrhonormal matrix $U=Q\tilde{U}$.
The stochastic SVD provides a robust means to perform large-scale matrix
decompositions, independently of the intrinisc structure of the data matrix.
However, this approach yields an approximate solution whose quality must be
assessed in the context of the original application.
In this work, the purpose of the stochastic SVD is to obtain the principal
components of a given dataset, i.e., to obtain the singular vectors to the
left of an SVD of a data matrix. We know that the obtained vectors form a
basis for a subspace, so that its precision can be assessed at each execution
of the stochastic SVD. This assessment was done using a distance between
subspaces as described in the next section.
#### 2.1.2 Distance between subspaces
Let $F$ and $G$ be subspaces of $\mathbbm{R}^{m}$, with
$p=\textsf{Dim}(F)\geq\textsf{Dim}(G)\geq q\geq 1$. The princpal angles,
$\theta_{1},\dots,\theta_{q}\in[0,\pi/2]$, between $F$ and $G$ are recursively
defined as being real the numbers $\theta_{k}$ such that
$\cos(\theta_{k})=\max_{u\in F,v\in G}|u^{\top}v|=|u_{k}^{\top}v_{k}|,$
subject to
$u^{\top}u=v^{\top}v=1,\qquad u^{\top}u_{i}=v^{\top}v_{i}=0\ ,\qquad
i=1,\dots,k-1.$
The vectors $\\{u_{1},\dots,u_{q}\\}$ e $\\{v_{1},\dots,v_{q}\\}$ are the so-
called principal vectors between the subspaces $F$ e $G$.
The greatest principal angle is related to the notion of distance between
subspaces of the same dimension (Golub & Van Loan,, 1996). Thus, if $p=q$ then
$dist(F,G)=\sqrt{1-\cos(\theta_{p})^{2}}$ $=\hskip
2.0pt\mathrm{sen}(\theta_{p})$. In a practical way, if the columns of
$Q_{F}\in\mathbbm{R}^{m\times p}$ and those of $Q_{G}\in\mathbbm{R}^{m\times
p}$ define orthonormal bases for $F$ and $G$, respectively, the cosine of the
greatest principal angle is determined by computing an SVD of the matrix
$Q_{F}^{\top}Q_{G}$ and taking the smallest of its singular values. It is
important to note that
$0\leq dist(F,G)\leq 1.$
The distance will be zero if $F=G$ and one if $F\cap G^{\perp}\neq\\{0\\}$,
with $G^{\perp}$ denoting the orthogonal complement of $G$, i.e., the space of
all orthogonal vectors to $G$.
### 2.2 Diffusion maps
We have explored a more recent non-linear technique for dimensionality
reduction and manifold learning, the so-called diffusion maps, which re-
organizes the data according to a reduced set of parameters related to the
approximate intrinsic geometry of the underlying phenomena (Lafon & Lee,,
2006; de la Porte et al.,, 2008). The reduced set of parametrs are computed
from the eigenvalues and eigenvectors of a diffusion operator on data. It is
robust to noise and outliers in data, and can be efficiently computed when its
application is properly designed.
Consider that the set of $n$ observations
$y_{1},\dots,y_{n}\in\mathbbm{R}^{p}$ approximately sample one or more non-
linear manifolds, each $y_{i}$ being a vector of all pixels of a video frame,
for instance. We describe the diffusion maps algorithm in four steps. The
first step consists in building an $n\times n$ matrix $W$ of pairwise (local)
similarities between the observations. The similarities are defined by a
kernel function $K:Y\times Y\rightarrow\mathbbm{R}$, satisfying
$K(y_{i},y_{j})=K(y_{j},y_{i})$ and $K(y_{i},y_{j})\geq 0$. In the present
work, we have used the heat kernel given as
$W_{ij}=K(y_{i},y_{j})=e^{-\frac{||y_{i}-y_{j}||}{\epsilon}^{2}}.$ (5)
By choosing the parameter $\epsilon$ one can adjust the size of the
neighborhood with which to compute similarities, based on prior knowledge of
the structure and density of data (de la Porte et al.,, 2008) and on sparsity
considerations. In this work $\epsilon$ was taken as the largest squared
Euclidean distance between all datapoints, although a significantly smaller
$\epsilon$ could have been chosen, leading to a sparser matrix $W$ thus
increasing the efficiency of the algorithm.
The second stage consists in constructing the diffusion matrix $P$, which is a
stochastic matrix, whose lines are normalized to 1. The matrix $P$ is obtained
by
$P=D^{-1}W,$ (6)
where $D$ is a diagonal matrix whose entries are given by
$D_{ii}=\sum_{j=1}^{n}W_{ij},\quad\text{for }i=1,\dots,n.$ (7)
Each entry of $P$ provides the pairwise connectivity of the data in an
underlying similarity graph. The graph can be seen as a Markov chain on the
data points whose transition matrix is $P$; each entry $P_{ij}$ represents the
probability of transitioning from data point $i$ to $j$ in one diffusion step.
In other words, this is the probability of clustering together data points $i$
and $j$ in one step. When taking powers of the transition matrix $P$, one
increases the number of steps taken to cluster nodes $i$ and $j$ to form a
manifold.
The third step consists in computing the spectral decomposition of the
transition matrix $P$, thus obtaining the eigenvalues $\lambda_{i}$ and
corresponding eigenvectors $\psi_{i}$. Since the matrix $P$ is stochastic, its
greatest eigenvalue in absolute value, $\lambda_{0}$, equals $1$. When the
transition matrix is positive definite, $P$ has a sequence of positive
eigenvalues sorted in decreasing order
$1=\lambda_{0}\geq\dots\geq\lambda_{n-1}>0.$ (8)
The last step of diffusion maps is to actually perform the dimensionality
reduction. This is done by discarding eigenvalues of smaller indices. In
selecting the largest $q$ eigenvalues, we obtain a new feature vector
$\widetilde{y_{j}}\in\mathbbm{R}^{q}$, using the diffusion map given by
$\widetilde{y_{j}}=\begin{bmatrix}\lambda_{1}^{t}\psi_{j1}\\\ \vdots\\\
\lambda_{q}^{t}\psi_{jq}\\\ \end{bmatrix}\ ,$ (9)
with, e.g., $\psi_{j1}$ being the $j-$th component of eigenvector $\psi_{1}$
and $t>0$ being a parameter corresponding to the power of the matrix $P$,
whose effect is to cluster through the Markov chain in $t$ steps.
## 3 Image-Based Stochastic Modeling
Our main goal is to devise a tractable stochastic model whose samples enable
simulating a video of the observed behavior of ocean surface patterns. The
same approach should be useful for image-based modeling and tracking of other
continuous-time deformations. In addition to provide a framework for machine
learning and inference, stochastic modeling is necessary to account for
aspects of the phenomenon that can be difficult or impossible to model
explicitly in an efficient manner.
For the automatic construction of the model from video, towards an application
of recognizing ocean patterns, we propose the use of Bayesian methods of
probabilistic inference. The use of these methods require a “training” or
“learning” stage from enough input data and then strategies for automatically
fitting of the stochastic model to newly observed data, which we call the
“testing” or “tracking” stage (Mumford & Desolneux,, 2010).
The stochastic model will be built from the data-driven model given by either
diffusion maps or PCA. Thus, consider that patterns of a signal are to be
modeled, where $S(t)\in\mathbbm{R}^{q}$ is the video frame at instant $t$ in
the diffusion map or PCA model space. In this work, we specifically reduce the
frames to 3 dimensions, thus $q\in 3$, i.e., the motion patterns in the video
are to be modulated by only 3 parameters.
Using the Bayesian inference paradigm, one seeks to infer the state of the
random variable $S(t)$ on a new time point given observable data $I(t)$ at
other time points. In this work some examples of observable variables are:
* •
Number of sample points.
* •
The local curvature between sample points.
* •
The variation of sample time.
In order to build the stochastic model, it is necessary to define a
probability function $P(S(t),$ $I(t))$. The inference of $S(t)$ is carried out
using the _a posteriori_ probability $P(S(t)|I(t))$ through Bayes’ rule:
$P(S(t)|I(t))=\dfrac{P(S(t))P(I(t)|S(t))}{P(I(t))},\text{ with
}P(I(t))=\sum_{S(t)}P(S(t))P(I(t)|S(t))$ (10)
The above general approach can lead to three problems (Mumford & Desolneux,,
2010):
* •
The construction of the probability model, $P$;
* •
To find an algorithm to maximize the _a posteriori_ probability;
* •
To optimize the parameters $I$ of the model as to optimally fit the data.
In order to validate the stochastic model, one must take samples to produce a
stochastic simulation, and test if these reproduce the observed real-world
flows of the sea surface. We have been actively working on adequate specific
techniques for the above three problems for the overall approach proposed
here.
## 4 Results and Discussion
The experiments were preformed in the Scilab free software language and
environment, together with the image processing toolbox SIP.
The data consisted of a video of a beach front with a resolution of $360\times
640$ and $21\min 06s$ duration, illustrated in Figure 2.111HD Florida Beaches
Sunset, Powerfloe Network, http://www.youtube.com/watch?v=0GBpGRXaruE,
02/10/2012. Used with permission. This video was chosen since it provides a
recording without camera movement, as it simplifies the modeling of the
problem without requiring a very large number of frames.
We extracted four $15s$ video clips to be analyzed from the original video,
forming the “training set” used to learn our model. For each clip all frames
of the video were extracted and converted to grayscale, resulting in a set of
$450$ images having $360\times 640=230\,400$ pixels or dimensions. These clips
sample the temporal evolution of the behaviour of ocean surface patterns. The
extracted frames of each clip go through dimensionality reduction using both
stochastic PCA and diffusion maps in order to perform a comparative study of
the power of each approach to represent the underlying patterns.
Figure 2: Sample image of the video used in the experiments. The larger goal
is to obtain a model of the 3D structure and motion of the ocean surface. We
begin to tackle this by the construction of an efficient stochastic model of
the image patterns in the video.
##### Reliability of the Stochastic PCA
As we have used PCA with a stochastic SVD algorithm, the obtained subspace at
each realization of the method has a certain variance. In order to assess
this, we computed the precision of the obtained subspaces for multiple runs of
stochastic PCA, using the distance betweeen subspaces described in Section
2.1.2.
To the best of our knowledge, there is no information in the literature about
the number of samples in stochastic SVD needed to obtain a set that generates
a subspace that provides a good enough approximation to the principal
components. For each set of images, the stochastic PCA was executed five
times, followed by the pairwise computation of subspace distances in order to
assess the variance of the result. Table 1 shows the average distance between
the obtained subspaces, as well as the standard deviation for each video clip.
Table 1: The average distance between the subspaces obtained through multiple runs of stochastic PCA, together with the standard deviation, for four different clips of the same original video. | Average distance | Standard deviation
---|---|---
Group 1 | 0.1099 | 0.003
Group 2 | 0.1105 | 0.039
Group 3 | 0.1064 | 0.056
Group 4 | 0.1092 | 0.005
Despite the average distances between the subspaces at each execution of
stochastic PCA being close to $0.1$, which is around $10\%$ of the maximum
distances between two observed subspaces, this result is acceptable for the
applications up to this point. With the low standard deviation values, we have
confirmed that the proposed algorithm gives consistent results in practice.
##### Diffusion Maps Vs. Stochastic PCA
The results of applying dimensionality reduction using stochastic PCA and
diffusion maps for the first videoclip are shown in Figure 3.
(a) Stochastic PCA (b) Diffusion maps
Figure 3: Projection of each frame of the first videoclip ($15\min$) in
$\mathbbm{R}^{3}$ using 3(a) one instance of stochastic PCA and 3(b) diffusion
map at scale $t=2$.
From Figure 3 it is not possible to be conclusive about which technique has
better potential as part of the proposed modeling pipeline. To further
investigage this, we compare the decay of the eigenvalues used in the
diffusion map with that of the eigenvalues used in the stochastic PCA,
relative to the largest eigenvalues. The diffusion map eigenvalues present a
much shaper decay than the eigenvalues used in the stochastic PCA, Figure 4.
This indicates that diffusion maps can represent the observed flow with fewer
parameters.
Figure 4: The decay of the first $100$ eigenvalues used in the diffusion map
compared to that of the eigenvalues used in the stochastic PCA, relative to
the largest eigenvalues, for the first videoclip.
## 5 Conclusion and Future work
The image-based fluid modeling approach proposed in this work is part of a
larger effort with applications to the online monitoring of ocean conditions
for aiding offshore processes of oil drilling, among others.
With dimensionality reduction it is possible to represent the original data
with a data-driven model with a small number of parameters. Using diffusion
maps, these parameters reveal intrinsic structure of the original data,
keeping the most relevant observable characteristics, while our experiments
show that PCA was not as efficient in this sense. Building on such reduced
model we proposed the main ideas of a stochastic model to account for
remaining factors.
Work is underway to produce techniques for detailed machine learning and
inference based on this model. The final complete model should be able to
capture the patterns of fluid surface motion in a way that, when statistically
sampling from this model, we should obtain a sequence of images that can
reproduce the observed fluid behavior through time. In the future this can
then be used to help infer the state of the ocean surface in a new video
sequence. We have also been considering the use of video streams from multiple
views to improve quality and robustness.
## References
* de la Porte et al., (2008) de la Porte, J., Herbst, B. M. Hereman, H. & van der Walt, S. J. An Introduction to Diffusion Maps. Proceedings of the 19th Symposium of the Pattern Recognition Association of South Africa; 26-28 de novembro de 2008; Prasa. Cape Town: F. Nicolls, University of Cape Town, South Africa, 2008. pp. 15-25.
* Golub & Van Loan, (1996) Golub, G. H. & Van Loan, C. F. Matrix computations. 3 ed. Baltimore : Johns Hopkins University Press, 1996.
* Halko et al., (2010) Halko, N.; Martinsson, P. G.; Shkolnisky, Y. & Tygert, M. An algorithm for the principal component analysis of large data sets. SIAM Journal on Scientific Computing, vol. 33, n. 5, pp. 2580-2594, 2011.
* Halko et al., (2011) Halko, N.; Martinsson, P. G. & Tropp, J. A. Finding structure with randomness: Probabilistic algorithms for constructing approximate matrix decompositions. SIAM Review, vol. 53, n. 2, pp. 217-288, 2011.
* Hastie & et al., (2013) Hastie, T.; Tibshirani, R. & Friedman, J. The elements of statistical learning: Data mining, inference, and prediction. Springer, 2013.
* Jolliffe, (2002) Jolliffe, I. T. Principal component analysis. 2 ed. New York : Springer, 2002.
* Lafon & Lee, (2006) Lafon, S. & Lee, A. Diffusion maps and coarse-graining: A unified framework for dimensionality reduction, graph partitioning, and data set parameterization. Pattern Analysis and Machine Intelligence, IEEE Transactions on, vol. 28, n. 9, pp. 1393-1403, 2006.
* Mumford & Desolneux, (2010) Mumford, D. & Desolneux, A. Pattern theory: the stochastic analysis of real-world signal. CRC Press, 2010.
|
arxiv-papers
| 2013-09-30T16:39:21 |
2024-09-04T02:49:51.718315
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Mauro de Amorim, Ricardo Fabbri, Lucia Maria dos Santos Pinto and\n Francisco Duarte Moura Neto",
"submitter": "Ricardo Fabbri",
"url": "https://arxiv.org/abs/1309.7912"
}
|
1309.7928
|
DPF2013-45 September 30, 2013
Measurements of $B\rightarrow DK$ decays to constrain the CKM Unitarity
Triangle angle $\gamma$ and related results at LHCb
Daniel Craik 111On behalf of the LHCb Collaboration.
Department of Physics
University of Warwick, Coventry, UK
> Constraints on the CKM angle $\gamma$ are presented from GLW, ADS, and GGSZ
> analyses of $B^{\pm}\rightarrow DK^{\pm}$ at the LHCb experiment. The
> branching fractions of $B^{0}\rightarrow\overline{D}^{0}K^{+}\pi^{-}$ and
> $B_{s}^{0}\rightarrow\overline{D}^{0}K^{-}\pi^{+}$ are also reported,
> measured relative to the related mode
> $B^{0}\rightarrow\overline{D}^{0}\pi^{+}\pi^{-}$.
> PRESENTED AT
>
>
>
>
> DPF 2013
> The Meeting of the American Physical Society
> Division of Particles and Fields
> Santa Cruz, California, August 13–17, 2013
>
## 1 Measurements of $\gamma$ from $B^{\pm}\rightarrow DK^{\pm}$
The CKM angle
$\gamma=\text{arg}(-V_{\text{ud}}V^{*}_{\text{ub}}/V_{\text{cd}}V^{*}_{\text{cb}})$
is currently the least well-constrained angle in the Unitarity Triangle. So
far, the most-sensitive measurements of $\gamma$ from a single experiment have
been performed by Belle [1] and BaBar [2]. These measurements yield values of
$\left(68^{+15}_{-14}\right)^{\circ}$ and
$\left(69^{+17}_{-16}\right)^{\circ}$, respectively.
Tree-level processes such as $B^{\pm}\rightarrow DK^{\pm}$ provide a
theoretically clean measurement of $\gamma$ with no contributions from new
physics processes. This measurement can be compared with measurements from
loop-mediated processes, which are sensitive to new physics, to provide a test
of the Standard Model. The current limits on the CKM Unitarity Triangle due to
tree-level and loop processes, as calculated by the CKMFitter group [3], are
shown in Fig. 1.
Figure 1: Constraints on the CKM Unitarity Triangle due to (left) tree-level
processes and (right) loop-mediated processes.
### 1.1 GLW/ADS analysis of $B^{\pm}\rightarrow DK^{\pm}$ and
$B^{\pm}\rightarrow D\pi^{\pm}$
The GLW method [4] uses $D$ decays to $C\\!P$ eigenstates such as $K^{+}K^{-}$
and $\pi^{+}\pi^{-}$. Decays can proceed either via a $D^{0}$ or a
$\overline{D}^{0}$ with a phase difference of $\delta_{B}+\gamma$. Suppression
in the decay via $D^{0}$ with respect to the $\overline{D}^{0}$ decay limits
interference to $\mathcal{O}(10\,\%)$ in $B^{\pm}\rightarrow DK^{\pm}$ and
$\mathcal{O}(1\,\%)$ in $B^{\pm}\rightarrow D\pi^{\pm}$.
The ADS method [5] uses $D$ decays to quasi-flavour-specific states such as
$\pi^{+}K^{-}$ and $\pi^{-}K^{+}\pi^{+}\pi^{-}$. Here the suppression of one
of the $B$ decays is partially balanced by the suppression of one of the $D$
decays, giving larger interference terms while also introducing an additional
phase shift of $\delta_{D}$.
Analyses have been performed on $B^{\pm}\rightarrow DK^{\pm}$ and
$B^{\pm}\rightarrow D\pi^{\pm}$ with the $D$ meson reconstructed from the
final states $K^{+}K^{-}$, $\pi^{+}\pi^{-}$, $K^{+}\pi^{-}$, $\pi^{+}K^{-}$,
$K^{-}\pi^{+}\pi^{+}\pi^{-}$ and $\pi^{-}K^{+}\pi^{+}\pi^{-}$ using LHCb data
corresponding to $1\,\text{fb}^{-1}$ of $pp$ collisions at a centre of mass
energy of 7 TeV [6, 7]. The invariant mass distributions of the two- and four-
body suppressed ADS modes are shown in Fig. 2 and Fig. 3, respectively.
Figure 2: Fits to the invariant mass distributions of the two-body suppressed
ADS mode $\pi^{\mp}K^{\pm}$ in (top) $B^{\mp}\rightarrow DK^{\mp}$ and
(bottom) $B^{\mp}\rightarrow D\pi^{\mp}$. The $B^{\mp}\rightarrow DK^{\mp}$
and $B^{\mp}\rightarrow D\pi^{\mp}$ components are shown in red and green,
respectively. The shaded component indicates partially reconstructed
backgrond, the dashed magenta line corresponds to partially reconstructed
$\Lambda_{b}^{0}\rightarrow\Lambda_{c}^{+}h^{-}$ and the total shape also
includes a combinatoric background. Figure 3: Fits to the invariant mass
distributions of the four-body suppressed ADS mode
$\pi^{\mp}K^{\pm}\pi^{\pm}\pi^{\mp}$ in (top) $B^{\mp}\rightarrow DK^{\mp}$
and (bottom) $B^{\mp}\rightarrow D\pi^{\mp}$. The $B^{\mp}\rightarrow
DK^{\mp}$ and $B^{\mp}\rightarrow D\pi^{\mp}$ components are shown in red and
green, respectively. The shaded component indicates partially reconstructed
backgrond, the dashed magenta line corresponds to partially reconstructed
$B_{s}^{0}\rightarrow DK^{-}\pi^{+}$ and the total shape also includes a
combinatoric background.
The observables measured are the ratio of $DK$ to $D\pi$ for each $D$ final
state,
$R_{K/\pi}^{f}=\frac{\Gamma\left(B^{-}\rightarrow D\left[\rightarrow
f\right]K^{-}\right)+\Gamma\left(B^{+}\rightarrow
D\left[\rightarrow\overline{f}\right]K^{+}\right)}{\Gamma\left(B^{-}\rightarrow
D\left[\rightarrow f\right]\pi^{-}\right)+\Gamma\left(B^{+}\rightarrow
D\left[\rightarrow\overline{f}\right]\pi^{+}\right)}\,,$
the charge asymmetry for each final state,
$A_{h}^{f}=\frac{\Gamma\left(B^{-}\rightarrow D\left[\rightarrow
f\right]h^{-}\right)-\Gamma\left(B^{+}\rightarrow
D\left[\rightarrow\overline{f}\right]h^{+}\right)}{\Gamma\left(B^{-}\rightarrow
D\left[\rightarrow f\right]h^{-}\right)+\Gamma\left(B^{+}\rightarrow
D\left[\rightarrow\overline{f}\right]h^{+}\right)}\,,$
and the ratio of the suppressed to favoured modes for $D\rightarrow K\pi$ and
$D\rightarrow K\pi\pi\pi$,
$R_{h}^{\pm}=\frac{B^{\pm}\rightarrow
D\left[f_{\text{sup}}\right]h^{\pm}}{B^{\pm}\rightarrow
D\left[f\right]h^{\pm}}\,.$
The values obtained for each of these observables can be found in Refs. [6,
7]. These variables serve as inputs for the combined $\gamma$ measurements in
Section 1.3 and Section 1.4.
### 1.2 GGSZ analysis of $B^{\pm}\rightarrow DK^{\pm}$
The GGSZ method [8] exploits the variation of the strong phase $\delta_{D}$
across the Dalitz plot in $D$ decays to three-body self-conjugate states such
as $K_{S}^{0}\pi^{+}\pi^{-}$ and $K_{S}^{0}K^{+}K^{-}$. The Dalitz plot is
divided into bins, as shown in Fig. 4, chosen to maximise statistical
sensitivity. The populations of $B^{+}$ and $B^{-}$ decays in each bin are
given by
$N^{+}_{\pm i}=h_{B^{+}}\left[K_{\mp i}+(x^{2}_{+}+y^{2}_{+})K_{\pm
i}+2\sqrt{K_{i}K_{-i}}(x_{+}c_{\pm i}\mp y_{+}s_{\pm i})\right]\,,$
$N^{-}_{\pm i}=h_{B^{-}}\left[K_{\pm i}+(x^{2}_{-}+y^{2}_{-})K_{\mp
i}+2\sqrt{K_{i}K_{-i}}(x_{-}c_{\pm i}\pm y_{-}s_{\pm i})\right]\,,$
where $K_{\pm i}$ is the efficiency corrected yield in bin $\pm i$ due to
$D^{0}$ flavour tagged events from BaBar [9, 10] and $c_{\pm i}$ and $s_{\pm
i}$ are the cosine and sine of the strong phase $\delta_{D}$ in bin $\pm i$
from CLEO-c [11].
Figure 4: Binning schemes used for the Dalitz plots of (left) $D\rightarrow
K_{S}^{0}\pi^{+}\pi^{-}$ and (right) $D\rightarrow K_{S}^{0}K^{+}K^{-}$. Bins
in the top-left half of the plots
($m^{2}_{K_{S}^{0}h^{-}}>m^{2}_{K_{S}^{0}h^{+}}$) are identified as $+i$ and
bins in the bottom-right half are labeled $-i$.
The remaining parameters are left free in the fit to the data: $h_{B^{\pm}}$
are normalisation factors for $B^{\pm}$, and
$x_{\pm}=r_{B}\text{cos}(\delta_{B}\pm\gamma)$ and
$y_{\pm}=r_{B}\text{sin}(\delta_{B}\pm\gamma)$ are the Cartesian parameters,
which are sensitive to $\gamma$.
Analyses have been performed on $B^{\pm}\rightarrow DK^{\pm}$ with the $D$
meson reconstructed in the final states $K_{S}^{0}\pi^{+}\pi^{-}$ and
$K_{S}^{0}K^{+}K^{-}$ using LHCb data corresponding to $1\,\text{fb}^{-1}$ of
$pp$ collisions at a centre of mass energy of 7 TeV [12] and
$2\,\text{fb}^{-1}$ of $pp$ collisions at a centre of mass energy of 8 TeV
[13]. The values obtained for the Cartesian parameters in the 8 TeV analysis
are
$x_{+}=\left(-8.7\pm 3.1\text{(stat.)}\pm 1.6\text{(syst.)}\pm
0.6\text{(ext.)}\right)\times 10^{-2}\,,$ $x_{-}=\left(\phantom{-}5.3\pm
3.2\text{(stat.)}\pm 0.9\text{(syst.)}\pm 0.9\text{(ext.)}\right)\times
10^{-2}\,,$ $y_{+}=\left(\phantom{-}0.1\pm 3.6\text{(stat.)}\pm
1.4\text{(syst.)}\pm 1.9\text{(ext.)}\right)\times 10^{-2}\,,$
$y_{-}=\left(\phantom{-}9.9\pm 3.6\text{(stat.)}\pm 2.2\text{(syst.)}\pm
1.6\text{(ext.)}\right)\times 10^{-2}\,,$
where the third uncertainty is due to the CLEO-c strong phase measurements
used in the fit.
Combining these values with the results from the 7 TeV analysis and fitting
for $\gamma$, $r_{B}$ and $\delta_{B}$ yields values of $(57\pm 16)^{\circ}$,
$(8.8^{+2.3}_{-2.4})\times 10^{-2}$ and $(124^{+15}_{-17})^{\circ}$,
respectively, where the values for $\gamma$ and $\delta_{B}$ are modulo
$180^{\circ}$. Two-dimensional projections of the confidence regions for these
parameters are shown in Fig. 5.
Figure 5: Two-dimensional projections of the confidence regions onto the
(left) $(\gamma,r_{B})$ and (right) $(\gamma,\delta_{B})$ planes. Contours
indicate the $1$, $2$ and $3\sigma$ boundaries and diamonds mark the central
values.
### 1.3 Combination of results from $1\,\text{fb}^{-1}$ measurements
The results in Section 1.1 and Section 1.2 are combined using a frequentist
approach to obtain a more constraining measurement of $\gamma$ [14]. In
addition to these results further measurements are included to improve the
fit: measurements of the strong phases and coherence factors for $D\rightarrow
K\pi$ and $D\rightarrow K\pi\pi\pi$ decays from CLEO-c[15], $C\\!P$ asymmetry
measurements of the neutral $D$ mesons from the Heavy Flavour Averaging
Group[16] and charm mixing parameters from LHCb[17]. A likelihood is
constructed from the measured observables as
${\cal
L}\left(\vec{\alpha}\right)={\displaystyle\prod_{i}}\xi_{i}\left(\vec{A}_{i}^{\text{obs}}|\vec{\alpha}\right)\,,$
where the sum is over the different measurements, $\vec{\alpha}$ is the set of
parameters and $\xi_{i}$ denotes the likelihood probability density functions
(PDFs) of the observables $\vec{A}_{i}^{\text{obs}}$. For most observables a
Gaussian PDF is assumed, however, where highly non-Gaussian behaviour is
observed, the experimental likelihood is used.
A combined $\gamma$ measurement has been performed including the results from
Section 1.1 and a subset of the results from Section 1.2 corresponding to
$1\,\text{fb}^{-1}$ of $pp$ collisions at a centre of mass energy of 7 TeV
[12]. The best-fit values and confidence intervals (modulo $180^{\circ}$) of
$\gamma$ are given in Table 1 and the $1-\text{CL}$ curves for $\gamma$ are
shown in Fig. 6.
combination | $\gamma$ | 68 % CL | 95 % CL
---|---|---|---
$DK$ | $72.0^{\circ}$ | $[56.4,86.7]^{\circ}$ | $[42.6,99.6]^{\circ}$
$D\pi$ | $18.9^{\circ}$ | $[7.4,99.2]^{\circ}\cup[167.9,176.4]^{\circ}$ | -
$DK$ and $D\pi$ | $72.6^{\circ}$ | $[55.4,82.3]^{\circ}$ | $[40.2,92.7]^{\circ}$
Table 1: Best-fit values and confidence intervals for $\gamma$ from the
combination of $DK$ and $D\pi$ measurements.
Figure 6: $1-\text{CL}$ curves for $\gamma$ from the combined
$1\,\text{fb}^{-1}$ GLW/ADS and $1\,\text{fb}^{-1}$ GGSZ measurements using
(left) only $DK$, (centre) only $D\pi$ and (right) both decay modes.
### 1.4 Combination including $3\,\text{fb}^{-1}$ GGSZ measurement
Another combination [18] has been performed that incorporates all of the
results reported in Section 1.2 but only those observables from Section 1.1
corresponding to $B^{\pm}\rightarrow DK^{\pm}$ decays. Mixing in the neutral
$D$ mesons is also neglected in the equations used for the observables in this
combination.
The best-fit values and confidence intervals (all modulo $180^{\circ}$) for
$\gamma$, $r_{B}$ and $\delta_{B}$ are given in Table 2. Figure 7 and Figure 8
show the $1-\text{CL}$ curve for $\gamma$, and the 2D projection of the
likelihood in $\gamma$ and $r_{B}$, respectively.
quantity | value | 68 % CL | 95 % CL
---|---|---|---
$\gamma$ | $67.2^{\circ}$ | $[55.1,79.1]^{\circ}$ | $[43.9,89.5]^{\circ}$
$r_{B}$ | $0.0923$ | $[0.0843,0.1001]$ | $[0.0762,0.1075]$
$\delta_{B}$ | $114.3^{\circ}$ | $[101.3,126.3]^{\circ}$ | $[88.7,136.3]^{\circ}$
Table 2: Best-fit values and confidence intervals for $\gamma$, $r_{B}$ and
$\delta_{B}$ from the combination of $DK$ measurements including GGSZ
measurements from $3\,\text{fb}^{-1}$ of data. Figure 7: $1-\text{CL}$ curve
for $\gamma$ from the combined $1\,\text{fb}^{-1}$ GLW/ADS and
$3\,\text{fb}^{-1}$ GGSZ measurements. Figure 8: Two-dimensional projection of
the confidence regions onto the $(\gamma,r_{B})$ plane. Contours show the $1$
and $2\sigma$ boundaries and markers indicate the central values.
## 2 Measurement of $B_{(s)}^{0}\rightarrow DK\pi$ branching fractions
The decay mode $B^{0}\rightarrow DK^{+}\pi^{-}$ has potential for a
significant future measurement of $\gamma$ [19, 20, 21]. Sensitivity to
$\gamma$ comes from the interference of $b\rightarrow c$ and $b\rightarrow u$
amplitudes of a similar magnitude. $B_{s}^{0}\rightarrow DK^{-}\pi^{+}$ and
the related mode $B_{s}^{0}\rightarrow D^{*}K^{-}\pi^{+}$ form important
backgrounds to this mode, therefore, an understanding of these modes is
necessary.
Branching fraction measurements of $B^{0}\rightarrow DK^{+}\pi^{-}$ and
$B_{s}^{0}\rightarrow DK^{-}\pi^{+}$, relative to the normalisation mode
$B^{0}\rightarrow D\pi^{+}\pi^{-}$, have been made using LHCb data
corresponding to $1\,\text{fb}^{-1}$ of $pp$ collisions at a centre of mass
energy of 7 TeV [22]. The invariant mass distributions of $D\pi\pi$ and
$DK\pi$ candidates where the $D$ is reconstructed from
$\overline{D}^{0}\rightarrow K^{+}\pi^{-}$ are shown in Fig. 9. The measured
relative branching fractions are
$\frac{{\cal B}\left(B^{0}\to\overline{D}^{0}K^{+}\pi^{-}\right)}{{\cal
B}\left(B^{0}\to\overline{D}^{0}\pi^{+}\pi^{-}\right)}=0.106\pm
0.007\,\text{(stat.)}\pm 0.008\,\text{(syst.)}\,,$ $\frac{{\cal
B}\left(B_{s}^{0}\to\overline{D}^{0}K^{-}\pi^{+}\right)}{{\cal
B}\left(B^{0}\to\overline{D}^{0}\pi^{+}\pi^{-}\right)}=1.18\pm
0.05\,\text{(stat.)}\pm 0.12\,\text{(syst.)}\,.$
Figure 9: Fits to the $B^{0}_{(s)}$ candidate invariant mass distributions for
the (a) $D\pi\pi$ and (b) $DK\pi$ samples. Data points are shown in black, the
full fitted PDFs as solid blue lines and the components as detailed in the
legends.
These relative measurements yield absolute branching fractions of
${\cal B}\left(B^{0}\to\overline{D}^{0}K^{+}\pi^{-}\right)=(9.0\pm
0.6\,\text{(stat.)}\pm 0.7\,\text{(syst.)}\pm 0.9({\cal B}))\times 10^{-5}\,,$
${\cal B}\left(B_{s}^{0}\to\overline{D}^{0}K^{-}\pi^{+}\right)=(1.00\pm
0.04\,\text{(stat.)}\pm 0.10\,\text{(syst.)}\pm 0.10({\cal B}))\times
10^{-3}\,,$
where the third uncertainty arises from the uncertainties on ${\cal
B}(B^{0}\to\overline{D}^{0}\pi^{+}\pi^{-})$. This is the most precise
measurement of ${\cal B}(B^{0}\to\overline{D}^{0}K^{+}\pi^{-})$ to date and
the first measurement of ${\cal B}(B_{s}^{0}\to\overline{D}^{0}K^{-}\pi^{+})$.
Although no quantitative analysis of the Dalitz plots has yet been attempted,
the Dalitz plot distributions obtained (corrected for efficiency) are
presented in Fig. 10.
Figure 10: Efficiency corrected Dalitz plot distributions for (a)
$B^{0}\rightarrow\overline{D}^{0}\pi^{+}\pi^{-}$, (b)
$B^{0}\rightarrow\overline{D}^{0}K^{+}\pi^{-}$ and (c)
$B_{s}^{0}\rightarrow\overline{D}^{0}K^{-}\pi^{+}$ candidates obtained from
the signal weights.
## 3 Conclusions and prospects
The $B^{\pm}\rightarrow DK^{\pm}$ decay mode offers an excellent opportunity
to measure the CKM angle $\gamma$ from Standard Model processes. The
combination in Section 1.4 gives the most sensitive measurement of $\gamma$
from a single experiment so far, yielding a value of $(67\pm 12)^{\circ}$.
This measurement is expected to improve further with the completion of a
GLW/ADS analysis on the remaining $2\,\text{fb}^{-1}$ of LHCb data currently
available. In addition, other modes such as $B^{0}\rightarrow DK^{+}\pi^{-}$
offer great prospects for future $\gamma$ measurements.
ACKNOWLEDGMENTS
This work is funded in part by the European Research Council under FP7 and by
the United Kingdom’s Science and Technology Facilities Council.
## References
* [1] K. Trabelsi [Belle Collaboration], arXiv:1301.2033 [hep-ex].
* [2] J. P. Lees et al. [BaBar Collaboration], Phys. Rev. D 87 (2013) 052015 [arXiv:1301.1029 [hep-ex]].
* [3] J. Charles et al. [CKMfitter Group Collaboration], Eur. Phys. J. C 41 (2005) 1 [hep-ph/0406184].
* [4] M. Gronau and D. Wyler, Phys. Lett. B 265 (1991) 172.
* [5] D. Atwood, I. Dunietz and A. Soni, Phys. Rev. Lett. 78 (1997) 3257 [hep-ph/9612433].
* [6] R. Aaij et al. [LHCb Collaboration], Phys. Lett. B 712 (2012) 203 [Erratum-ibid. B 713 (2012) 351] [arXiv:1203.3662 [hep-ex]].
* [7] R. Aaij et al. [LHCb Collaboration], Phys. Lett. B 723 (2013) 44 [arXiv:1303.4646 [hep-ex]].
* [8] A. Giri, Y. Grossman, A. Soffer and J. Zupan, Phys. Rev. D 68 (2003) 054018 [hep-ph/0303187].
* [9] B. Aubert et al. [BaBar Collaboration], Phys. Rev. D 78 (2008) 034023 [arXiv:0804.2089 [hep-ex]].
* [10] P. del Amo Sanchez et al. [BaBar Collaboration], Phys. Rev. Lett. 105 (2010) 121801 [arXiv:1005.1096 [hep-ex]].
* [11] J. Libby et al. [CLEO Collaboration], Phys. Rev. D 82 (2010) 112006 [arXiv:1010.2817 [hep-ex]].
* [12] R. Aaij et al. [LHCb Collaboration], Phys. Lett. B 718 (2012) 43 [arXiv:1209.5869 [hep-ex]].
* [13] R. Aaij et al. [LHCb Collaboration], LHCb-CONF-2013-004 (2013).
* [14] R. Aaij et al. [LHCb Collaboration], arXiv:1305.2050 [hep-ex].
* [15] N. Lowrey et al. [CLEO Collaboration], Phys. Rev. D 80 (2009) 031105 [arXiv:0903.4853 [hep-ex]].
* [16] Y. Amhis et al. [Heavy Flavor Averaging Group Collaboration], arXiv:1207.1158 [hep-ex].
* [17] R. Aaij et al. [LHCb Collaboration], Phys. Rev. Lett. 110 (2013) 101802 [arXiv:1211.1230 [hep-ex]].
* [18] R. Aaij et al. [LHCb Collaboration], LHCb-CONF-2013-006 (2013).
* [19] M. Gronau, Phys. Lett. B 557 (2003) 198 [hep-ph/0211282].
* [20] T. Gershon, Phys. Rev. D 79 (2009) 051301 [arXiv:0810.2706 [hep-ph]].
* [21] T. Gershon and M. Williams, Phys. Rev. D 80 (2009) 092002 [arXiv:0909.1495 [hep-ph]].
* [22] R. Aaij et al. [LHCb Collaboration], Phys. Rev. D 87 (2013) 112009 [arXiv:1304.6317 [hep-ex]].
|
arxiv-papers
| 2013-09-30T17:29:46 |
2024-09-04T02:49:51.725944
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Daniel Craik",
"submitter": "Daniel Craik",
"url": "https://arxiv.org/abs/1309.7928"
}
|
1309.7938
|
# The NIFFTE project
J. Ruza Corresponding author. D.M Asnerb R.G Bakerc J. Bundgaardd E. Burgette
M. Cunninghama J. Deavene D.L. Duked U. Greifed S. Grimesf M. Heffnera T.
Hillg D. Isenhowerh J.L. Klayc V. Kleinrathe N. Kornilovf A.B. Laptevi W.
Lovelandj T.N. Masseyf R. Meharchandi H. Quh S. Sangiorgioa B. Seilhana L.
Snydera S. Staveb G. Tatishvilib R.T. Thorntonh F. Tovessoni D. Towellh R.S.
Towellh S. Watsonh B. Wendte and L. Woodb
aLawrence Livermore National Laboratory Livermore CA 94550 USA
bPacific Northwest National Laboratory Richland WA 99354 USA
cCalifornia Polytechnic State University San Luis Obispo CA 93407 USA
dColorado School of Mines Golden CO 80401 USA
eIdaho State University Pocatello ID 83209 USA
fOhio University Athens OH 45701 USA
gIdaho National Laboratory Idaho Falls ID 83415 USA
hAbilene Christian University Abilene TX 79699 USA
iLos Alamos National Laboratory Los Alamos NM 87545 USA
jOregon State University Corvallis OR 97331 USA
E-mail [email protected]
###### Abstract
The Neutron Induced Fission Fragment Tracking Experiment (NIFFTE) is a double-
sided Time Projection Chamber (TPC) with micromegas readout designed to
measure the energy-dependent neutron-induced fission cross sections of the
major and minor actinides with unprecedented accuracy. The NIFFTE project
addresses the challenge of minimizing major sources of systematic
uncertainties from previous fission chamber measurements such as: target and
beam non-uniformities, misidentification of alpha and light charged particles
as fission fragments, and uncertainties inherent to the reference standards
used. In-beam tests of the NIFFTE TPC at the Los Alamos Neutron Science Center
(LANSCE) started in 2010 and have continued in 2011, 2012 and 2013. An
overview of the NIFFTE TPC status and performance at LANSCE will be presented.
###### keywords:
Fission; TPC; Cross Section
## 1 Introduction
Sensitivity studies indicate that high-accuracy measurements of fission cross
sections are needed and the Neutron Induced Fission Fragment Tracking
Experiment (NIFFTE) collaboration is developing a Time Projection Chamber
(TPC) to meet this need. Neutron-induced fission cross sections have been
measured with fission chambers (ICs) for decades, typically by placing a thin
target of the fissile material of interest into a neutron beam and detecting
the ionization left in a gas by the resulting fission products. To date,
fission chambers have been limited to measuring the energy deposited in a gas
by charged particles, and were therefore unable to distinguish between
different particles of the same energy. Fragments are subject to energy loss
and straggling within the target111A non negligible fraction of fission
fragments might become indistinguishable from lighter charged particles within
the fission chamber due to energy loss and straggling. resulting in a
systematic error for the measured cross section.
Based on well-established technology originally developed for high-energy
physics experiments, the NIFFTE TPC is uniquely suited to achieve
unprecedented accuracy. The detection principle of a TPC is rather simple:
particles ionizing the gas traverse a uniform electric field. The electrons
produced by the ionization are drifted by the electric field and collected at
a segmented readout plane. A two-dimensional projection of the particle’s
track is obtained by interpreting the amount of charge deposited in each
segment. Combining this two dimensional projection with the drift time of the
electrons provides a full, three-dimensional representation of the particle’s
path through the detector.
A schematic of a neutron-induced fission in the TPC is shown in figure 1. The
fissionable sample to be studied is placed in the center of the detector at
the midpoint of the cathode. The field cage covers the $\rm{54\,mm}$ drift
distance at each side of the sample and it is placed inside a pressure vessel
of $\rm{15\,cm}$ diameter that is filled up with a suitable gas mixture e.g.
Argon($90\,\%$)-Methane($10\,\%$) at a nominal $\rm{550\,torr}$. At each end
of the chamber charge is amplified and read out on a MICROMEGAS “pad” plane.
The signals are digitized and processed by custom-made EtherDAQ card
assemblies [1, 2] sampling at $50\,\rm{MHz}$ together with a software
framework developed for the data acquisition, analysis, and simulation tasks
[3]. An article describing the TPC hardware is in preparation.
Figure 1: Schematic of a neutron-induced fission in the target of the NIFFTE
TPC (left) and picture of the TPC installed in the LANSCE-WNR facility during
the 2012 run cycle.
## 2 Beam experiments
Building a TPC suitable for fission research is an ambitious task, involving
complex systems and state-of-the-art technology. As new hardware and software
capabilities are implemented, frequent evaluation is required. In-beam testing
takes place at the 4FP90L flight path of Los Alamos Neutron Science Center
(LANSCE) Weapons Neutron Research (WNR) facility [4] that makes use of an
unmoderate tungsten spallation target to produce a high-energy neutron
spectrum with energies ranging from $100\,\rm{keV}$ to $600\,\rm{MeV}$.
The first TPC experiments at LANSCE-WNR took place in 2010. Data were
collected first with 64 then 192 individual channels ($\rm{1\,\%}$ and
$\rm{3\,\%}$ of the full TPC), with blank carbon samples and on carbon
backings with $\rm{{}^{238}U}$ deposited onto them [5]. These early
measurements validated the TPC technique, as particle tracks could be
reconstructed for both fission fragments and alpha particles. During the 2011
LANSCE run cycle, one complete sextant (496 channels, $\rm{8\,\%}$ of the full
TPC) was instrumented, and data were collected with samples containing
$\rm{{}^{238}U}$, $\rm{{}^{235}U}$, and $\rm{{}^{239}Pu}$. During the most
recent run cycle (2012), the TPC was instrumented with one complete pad plane
(2976 channels, 1/2 of the full detector), and notably, for the first time in
the TPC project, information about the neutron time-of-flight was available.
## 3 Performance
Preliminary analysis of the data taken during the 2010 and 2011 LANSCE run
cycles have been reported [6, 7] and the 2012 run cycle data is currently
being analyzed. Data were collected on three samples: a split sample
(semicircular deposits) containing $\rm{{}^{238}U}$ and $\rm{{}^{235}U}$; a
split sample with $\rm{{}^{235}U}$ and $\rm{{}^{239}Pu}$; and a sample with
$\rm{{}^{239}Pu}$ electro-plated onto metallized polyproplylene. In adition,
off-beam data was collected with the use of $\rm{{}^{252}Cf}$ and
$\rm{{}^{244}Cm}$ bottom sources to fine tune the acquisition software and the
analysis codes. All results included in this paper correspond to the 2012 data
taking period and should be considered preliminary.
### 3.1 Neutron Time of Flight
Neutron-induced fission depends on the energy of the incident neutron. The
neutron Time-Of-Flight (nTOF) between the spallation target at WNR and the
fissionable target placed inside the TPC can be folded together with the
overall WNR beam spacing structure -nominal $1.8\,\rm{\mu s}$ pulse spacing
between neutron bunches -to reproduce the energy spectrum of the neutrons
inducing fission events in the TPC. Figure 2 shows a reconstruction of the
nTOF for the NIFFTE TPC during the 2012 run cycle. The peak around
$28\,\rm{nsec}$ corresponds to photo-induced fission in the TPC target and,
can be used to obtain an absolute calibration for both, the distance of the
TPC relative to the spallation target and the energy of the incoming neutrons.
The width of the photo-fission peak serves as a tool to determine the
systematic error on the energy of the incident neutron. During 2012, the
NIFFTE TPC was able to achieve $2.5-3.5\,\rm{nsec}$ FWHM for the photo-fission
peak, which is competitive with the resolution obtained by previous fission
chamber experiments.
Figure 2: Neutron time of flight reconstruction for neutron-induced fission
events in the TPC and zoom of the observed photo-fission peak. A split sample
(semicircle deposits) containing $\rm{{}^{238}U}$ and $\rm{{}^{235}U}$ was
used during this measurement.
### 3.2 Tracking and Particle Identification
One of the major advantages of using a TPC to study neutron-induced fission
cross sections is the ability to reconstruct and study the features of the
ionizing tracks that the different particles generate on their way through the
active volume of NIFFTE.
One example of how particle identification (distinguishing alpha and light
charged particles from fission fragments) might be done using the tracking
information of the TPC can be shown by comparing the length of the particle’s
track to its total energy deposited in the TPC (see figure 3). Fission
fragments, have relatively short tracks and deposit large amount of energy in
the fill gas, while light particles (alpha particles) have longer tracks and
deposit less energy.
Figure 3: Plot of length versus energy of track reconstructions in the TPC for
a $\rm{{}^{252}Cf}$ buttom source. On the plot there are three regions that
can be easily spotted, an almost vertical line at low energies with lengths
from $\rm{1\,cm}$ up to $\rm{6\,cm}$ containing $\alpha$-particles and
protons, a big elliptical zone centered around $\rm{80\,MeV}$ with tracks from
$\rm{1.5\,cm}$ up to $\rm{3\,cm}$ gathering most of the fission fragment
population and a small shadow to the elipse with track lengths around
$\rm{1\,cm}$ and energies between $\rm{5}$ and $\rm{40\,MeV}$. The shadow
region has been succesfully identified as fission fragments that come out of
the target at high polar angles, suffering from straggling, and hit on the
concentric copper ring of $\rm{2}$ by $\rm{4\,cm}$ diameter that holds the
polypropylene backing foil where the $\rm{2\,mm}$ diameter $\rm{{}^{252}Cf}$
spot was deposited.
Using precise tracking information, e.g. the differential ionization pattern
of the particle within the gas, allows the search of the specific Bragg peaks
created by the particles and the comparison of their ionization pattern with
simulations and theoretical models to provide an ultimate tool for particle
identification. Fission fragments deposit most of the energy at the begining
of their short track through the TPC, while $\alpha$ particles and protons
penetrate the gas further and have their Bragg peak towards the end of their
path through the TPC (see figures 4 and 5).
Figure 4: Ionization pattern of a fission fragment in the NIFFTE TPC. The
Bragg peak is located at the begining of the track for these particles.
Finally, tracking information also enables tracing of each individual particle
back to a well-defined origin in the sample, which ultimately allows for a
minimization on the systematics errors such as beam spot size, and sample
uniformity. This makes it possible to study two or more actinides on the same
backing, as described in [5], and to check sample uniformity using the decay
activity of the actinide deposits.
Figure 5: Ionization pattern of an observed $\alpha$-track in the NIFFTE TPC
(black) and a theoretical prediction that includes the ionization nature of an
$\alpha$-like particle together with the specific diffusion expected in the
TPC for this particular track (blue). In red the residual of observed and
modeled ionization pattern is shown. Figure 6: Track vertex reconstructions of
$\alpha$-like events emerging from a $\rm{2\,mm}$ spot source of
$\rm{{}^{252}Cf}$ placed in the center of the NIFFTE TPC. In the figure a
picture of the sample used can be seen.
Figure 6 shows the vertex reconstructions for spontaneous $\alpha$-like events
from a $\rm{{}^{252}Cf}$ bottom source. The tracks seem to originate from the
center of the cathode plane, coming out of a circle of approximately
$\rm{2\,mm}$ diameter, which is consistent with the actual dimension of the
$\rm{{}^{252}Cf}$ deposit. Note that the two different backing materials can
easily be distinguished: one external ring of copper with inner and outer
diameters of $\rm{2\,cm}$ and $\rm{4\,cm}$ respectively, and a $\rm{2\,cm}$
diameter polypropylene foil on top of which the $\rm{{}^{252}Cf}$ was
deposited. This figure clearly shows the radiographic capabilities of NIFFTE
to map sample uniformities.
## 4 Conclusion
The NIFFTE TPC is being tested off-beam and in-beam at the neutron facility
LANSCE-WNR. Preliminary results (with half detector instrumented) indicate
that the TPC performance is meeting expectations. As the NIFFTE project moves
from development into production phase, the collaboration is optmistic that
the project is on-track to deliver neutron-induced fission cross sections with
unprecedented accuracy.
###### Acknowledgements.
This work was performed under the auspices of the U.S. Department of Energy by
Lawrence Livermore National Laboratory under contract No. DE-AC52-07NA27344.
## References
* [1] M. Heffner, V. Riot, and L. Fabris _IEEE Transactions on Nuclear Science_ 60, _3_ (June 2013).
* [2] L. Wood _et al._ , _Proceedings of the International Conference on Nuclear Data for Science and Technology_ , _Nuclear Data Sheets_ (2013).
* [3] S. Stave _et al._ , _Proceedings of the International Conference on Nuclear Data for Science and Technology_ , _Nuclear Data Sheets_ (2013).
* [4] P.W. Lisowski and K.F. Schoenberg, _Nuclear Instruments and Methods in Physics Search Section A_ 562, _910_ (2006).
* [5] W. Loveland, L. Yao _et al._ , _Proceedings of the International Conference on Nuclear Data for Science and Technology_ , _Nuclear Data Sheets_ (2013).
* [6] F. Tovesson, A.B. Laptev, R. Meharchand, _DOE-NE Technical Report LA-UR 11-05439_ , September 2011.
* [7] R. Meharchand, T. Hill, L. Snyder, _Los Alamos National Laboratory Report LA-UR 12-25452_ , October 2012.
|
arxiv-papers
| 2013-09-30T17:39:47 |
2024-09-04T02:49:51.732312
|
{
"license": "Public Domain",
"authors": "Ruz J., Asner D.M., Baker R.G., Bundgaard J., Burgett E., Cunningham\n M., Deaven J., Duke D.L., Greife U., Grimes S., Heffner M., Hill T.,\n Isenhower D., Klay J.L., Kleinrath V., Kornilov N., Laptev A.B., Loveland W.,\n Masseyf T.N., Meharchand R., Qu H., Sangiorgio S., Seilhan B., Snyder L.,\n Stave S., Tatishvili G., Thornton R.T., Tovesson F., Towell D., Towell R.S.,\n Watson S., Wendt B. and Wood L",
"submitter": "Jaime Ruz Armendariz",
"url": "https://arxiv.org/abs/1309.7938"
}
|
1309.7961
|
# Prospects For Precision Measurements with Reactor Antineutrinos at Daya Bay
Daya Bay Collaboration
(September 27, 2013)
###### Abstract
In 2012 the Daya Bay experiment made an unambiguous observation of reactor
antineutrino disappearance over kilometer-long baselines and determined that
the neutrino mixing angle $\theta_{13}$ is non-zero. The measurements of Daya
Bay have provided the most precise determination of $\theta_{13}$ to date.
This whitepaper outlines the prospects for precision studies of reactor
antineutrinos at Daya Bay in the coming years. This includes precision
measurements of sin22$\theta_{13}$ and $\Delta m^{2}_{ee}$ to $<$3%, high-
statistics measurement of reactor flux and spectrum, and non-standard physics
searches.
## I Introduction
Since the discovery of a non-zero value for $\theta_{13}$ DYB1 , Daya Bay has
provided the most precise measurement of sin22$\theta_{13}$ of $0.089\pm
0.010\pm 0.005$ with 139 days of 6 antineutrino detector (AD) data DYB2 in
June 2012. Recently, with 217 days of 6-AD data, we have obtained the most
precise measurement of sin22$\theta_{13}=0.090^{+0.008}_{-0.009}$ and the
first measurement of $\Delta m^{2}_{ee}=(2.59^{+0.19}_{-0.20})\times
10^{-3}~{}eV^{2}$.
Two more ADs were installed and have been in operation since October 19, 2012.
We present the expected precision in sin22$\theta_{13}$ and $\Delta
m^{2}_{ee}$ as a function of time in Sections II and III respectively. The
prospects for precise measurement of the reactor antineutrino flux and
spectrum are summarized in Section IV.
## II Precision Measurement of sin22$\theta_{13}$
Daya Bay’s currently published measurements of sin22$\theta_{13}$ are still
statistics limited. With the data collected with 8 ADs as of October 2013 the
statistical error will reach the current systematic uncertainty of 0.005 DYB2
. We have assessed improvements in the systematic uncertainty based on
calibration data and increased precision that is statistical in nature (such
as the uncertainties in estimating the 9Li/8He and fast neutron backgrounds).
Based on the projected systematic uncertainty we estimate that the total
uncertainty in sin22$\theta_{13}$ can be reduced from 0.008 to 0.006 in one
year and 0.003 after 3 to 4 years as shown in Fig. 1.
Figure 1: Projected total uncertainty in sin22$\theta_{13}$ for the Daya Bay
experiment. The current systematic uncertainty of 0.005 is used at the
beginning of the 8AD run. Figure 2: Significance with which CP violation can
be observed, by NOvA+T2K (left) and NOvA+T2K+LBNE (right), as a function of
the value of $\delta_{CP}$. Observation of CP violation is equivalent to
measuring $\delta_{CP}\neq 0,\pi$. The significance is calculated by
minimizing over normal and inverted hierarchies, as the hierarchy is assumed
to be unknown. The impact of $\theta_{13}$ precision is shown.
Daya Bay’s short-baseline antineutrino disappearance measurement is the most
precise reactor measurement of sin22$\theta_{13}$ and will likely remain the
most precise measurement of this fundamental parameter for the foreseeable
future. Precise measurement of this quantity may shed light on symmetries
between quarks and leptons at a fundamental level, and test unitarity of the
neutrino mixing matrix. Comparison of future long baseline accelerator
measurements of $\theta_{13}$ to Daya Bay will allow precision tests of the
neutrino-standard model interpretation of neutrino oscillations with
sensitivity to non-standard neutrino interactions and sterile neutrino
scenarios. The expected improvement in sin22$\theta_{13}$ by Daya Bay will
make $\theta_{13}$ the most precisely measured neutrino mixing angle and
extend the CP reach (shown in Fig. 2) and $\theta_{23}$ octant determination
of long baseline accelerator experiments.
## III Precision Measurement of $\Delta m^{2}_{ee}$
Our discovery of a large value for $\theta_{13}$ enables a measurement of the
effective mass splitting $\Delta m^{2}_{ee}$ with a precision comparable to
the MINOS measurement of $\Delta m^{2}_{\mu\mu}$. The disappearance of
electron antineutrinos over km-long baselines observed at Daya Bay is a
combination of antineutrino oscillations with mass splittings $\Delta
m^{2}_{31}$ and $\Delta m^{2}_{32}$.
Complementary to $\Delta m^{2}_{\mu\mu}$ determined by accelerator-based
experiments, Daya Bay will provide a precise measurement of $\Delta
m^{2}_{ee}$, as shown in Fig. 3, better than 3% after five years of data
collection with the 8-AD configuration. In addition to the neutrino mixing
angles and the Dirac CP-violating phase, the mass-squared differences are
crucial for understanding the nature of neutrinos. First, independent
determination of the three mass-squared differences will validate the sum rule
of neutrino mixing, $\Delta m^{2}_{21}+\Delta m^{2}_{32}+\Delta m^{2}_{13}=0$.
Any deviation from this sum rule could signal the existence of neutrinos
beyond three generations. In addition, comparison of $\Delta m^{2}_{ee}$ and
$\Delta m^{2}_{\mu\mu}$ will provide direct constraint to the neutrino mass
hierarchy (sign of $\Delta m^{2}_{31}$).
Figure 3: Expected Daya Bay uncertainty in $\Delta m^{2}_{ee}$ as a function
of running time with 8 ADs. The current systematic uncertainty is used at the
beginning of the run. The present error of the 6-AD data set is $0.2\times
10^{-3}~{}eV^{2}$.
Since the fundamental parameters $\Delta m^{2}_{32}$ and $\Delta m^{2}_{31}$
are directly related to the effective mass splitting $\Delta m^{2}_{ee}$ in
the L/E region for Daya Bay, the impact of the Daya Bay measurements in both
sin22$\theta_{13}$ and $\Delta m^{2}_{ee}$ with three years of 8-AD data is
illustrated in Fig. 4. Only statistical uncertainties are shown. The relative
systematic uncertainty, which is relevant for the oscillation analysis, is
still smaller than the expected statistical uncertainty of the far site. The
overall normalization is not constrained in the fit, and the best-fit
normalization is used in Fig. 4.
Figure 4: Expected L/E distribution of Daya Bay with 8 ADs in three
experimental halls, EH1, EH2, and EH3 after three years of running. Only
statistical uncertainties are shown.
## IV Precision Measurement of the Reactor Antineutrino Flux and Spectrum
Daya Bay collects reactor antineutrino data at a tremendous rate which enables
a precision measurement of the reactor antineutrino spectra in the four near
site detectors. For the data sets published to date based on six detectors,
Daya Bay has collected $>$300k antineutrino events. Two additional detectors,
one at the Ling Ao near site and one at the far site, have been operational
since October 2012. The Daya Bay experimental configuration allows spectral
and flux measurements as close as 360m from the reactor cores at the Daya Bay
site and 480m at the Ling Ao site. At these detector locations the measured
reactor antineutrino flux from the nearest reactor cores remains largely non-
oscillated in the standard 3-neutrino oscillation framework. The contribution
from the more distant reactors at $\sim$900m is approximately 20% (7.5%) of
the total event rate at Daya Bay (Ling Ao). The oscillation effects of
antineutrinos from the far reactor and resulting distortions in the measured
antineutrino spectrum can be corrected for in any spectral analysis assuming a
3-neutrino framework. With this unique configuration of multiple, functionally
identical antineutrino detectors at various baselines Daya Bay will be able to
report the following measurements and physics analyses:
### IV.1 High-precision measurement of the reactor antineutrino spectrum
Using data from the near and far sites Daya Bay will make a measurement of the
reactor antineutrino spectrum with high precision. Amongst all running reactor
experiments Daya Bay will collect the largest sample of antineutrino events
and achieve a $<$1% statistical uncertainty in a 2-year run over a large range
of energies at the near sites.
### IV.2 Test of the reactor antineutrino spectrum vs predictions and search
for new interactions
Using known reactor operation data such as thermal power output and fission
fraction evolution from reactor core simulations Daya Bay can predict the
expected non-oscillated energy spectrum of reactor antineutrinos emitted from
each reactor. The Daya Bay collaboration works closely with the reactor
company and relevant operation data are provided to the collaboration on a
regular basis. A precise comparison of the predicted reactor antineutrino
spectrum with the spectrum measured at Daya Bay will test our understanding
and calculations of antineutrino emitted from reactors. This is particularly
relevant in the context of recent discussions of our understanding of reactor
flux calculations. Discrepancies in the spectral shape may point to (a)
missing nuclear physics in the reactor spectrum predictions or (b) new physics
beyond the 3-neutrino framework including non-standard interaction (NSI)
effects, and sterile neutrinos. This spectral shape test is independent of the
absolute flux normalization and the uncertainties in the predicted total
antineutrino rate at Daya Bay. Due to the high statistics of the Daya Bay
measurement, the statistical uncertainty in the 2011-2012 Daya Bay data set is
already below the flux conversion uncertainty on the spectrum. Daya Bay’s
measurement of the reactor antineutrino spectrum will be ultimately limited by
our understanding of energy scale uncertainties and detector effects. Detector
studies and simulations are ongoing to improve the energy response model of
the Daya Bay detectors.
### IV.3 Absolute measurement of the reactor flux
In addition to a measurement of the spectral shape, Daya Bay will measure the
absolute reactor flux from the six Daya Bay and Ling Ao reactor cores. The
absolute flux measurement tests our understanding of reactor flux models
within the theoretical uncertainties of the predictions and the experimental
uncertainties associated with the absolute detection efficiencies of the Daya
Bay detectors. An absolute flux measurement in the antineutrino detectors at
Daya Bay will provide unique data points at the baselines of the Daya Bay
experiment (360-2,000 m) and will further our understanding if there is a
deficit in the measured reactor neutrino flux at short baselines, also known
as the “reactor anomaly”. An analysis of past reactor experiments compared
with predictions has revealed a discrepancy of about 5.7% in the absolute
antineutrino flux. While Daya Bay has demonstrated superb relative detector
uncertainties, an absolute measurement will be systematics limited. A
statistical precision of 0.1% will be achievable in the Daya Bay data set.
Improvements in the analysis may eventually reduce absolute detection
uncertainties to $\sim$1%. An absolute reactor flux measurement can test the
theoretical flux prediction with the uncertainty currently estimated at about
2.7%. We can compare the Daya Bay’s measurements to previous reactor flux
measurements, for example, the absolute Bugey-4 measurement with an
uncertainty of 1.4%. Daya Bay’s measurement of the absolute flux and reactor
antineutrino spectrum will provide important input to our understanding of the
“reactor anomaly”.
### IV.4 Study of the time-evolution of the reactor antineutrino flux
The large reactor antineutrino event rate measured and the unique combination
of baselines and reactor cores at Daya Bay allows a detailed study of the time
variation of the reactor antineutrino flux. This contains information on the
operation of the reactors as well as the evolution and isotopic composition of
the core’s fuel. Correlating the measured antineutrino flux with reactor
operations is of interest to reactor monitoring, the safeguard community, and
applied neutrino science. With six reactors and 4 near-site detectors Daya Bay
will provide the largest data set on reactor flux variations as a function of
time.
## References
* (1) F. P. An et al., Phys. Rev. Lett. 108, 171803 (2012).
* (2) F. P. An et al., Chin. Phys. C37, 011001 (2013).
|
arxiv-papers
| 2013-09-30T19:01:14 |
2024-09-04T02:49:51.738504
|
{
"license": "Public Domain",
"authors": "The Daya Bay Collaboration",
"submitter": "Kam Biu Luk",
"url": "https://arxiv.org/abs/1309.7961"
}
|
1310.0003
|
# Quantum criticality of geometric phase in coupled optical cavity arrays
under linear quench
Sujit Sarkar Poornaprajna Institute of Scientific Research, 4 Sadashivanagar,
Bangalore 5600 80, India.
e-mail: [email protected]
phone: 0091-80-23611836
Fax: 0091-80-23619034
###### Abstract
The atoms trapped in microcavities and interacting through the exchange of
virtual photons can be modeled as an anisotropic Heisenberg spin-1/2 lattice.
We study the dynamics of the geometric phase of this system under the linear
quenching process of laser field detuning which shows the XX criticality of
the geometric phase in presence of single Rabi frequency oscillation. We also
study the quantum criticality for different quenching rate in the presence of
single or two Rabi frequencies oscillation in the system.
PACS: 42.50.Pq, 03.65.Vf, 42.50.-p
Keywords: Cavity Quantum Electrodynamics, Geometric Phase and Quantum Optics.
Introduction: The recent experimental success in engineering strong
interaction between the photons and atoms in high quality micro-cavities opens
up the possibility to use light matter system as quantum simulators for many
body physics [1-16]. The authors of Ref. (hart1 ),(hart2 ) and (sujop ) have
shown that effective spin lattice can be generated with individual atom in the
micro-cavities that are coupled to each other via exchange of virtual photons.
The two states of spin polarization are represented by the two long lived
atomic levels in the system.
A Many body Hamiltonians can be created and probed in coupled cavity arrays.
In our previous study, we have explained explicitly the basic physics of the
formation of micro-optical cavity sujop . The atoms in the cavity are used for
detection and also for generation of interaction between photons in the same
cavity. As the distance between the adjacent cavities is considerably larger
than the optical wave length of the resonant mode, individual cavities can be
addressed. This artificial system can act as a quantum simulator. In this
optical cavities system one can find different quantum phases of polariton ( a
combined excitations of atom-photon interactions. ) by using the spin model
that conserve the total number of excitations. This micro-cavity system shows
the different quantum phases and quantum phase transitions.
Quantum Phase Transition (QPT) associate with the fundamental changes that
occurs in the macroscopic nature of the matter at zero temperature due to the
variation of some external parameter. The quantum phase transitions are
characterized by the drastic change in the ground state properties of the
system driven by the quantum fluctuations ss .
Motivation of this research paper:
To the best of our knowledge, here we not only study the dynamics of geometric
phase but also solve the nature of criticality explicitly under laser field
detuning quenching process for different quenching rate which is absent in the
previous literature of cavity QED [1-16].
Here we mention very briefly the essence of the geometric phase in the
condensed matter. The geometric phases have been associated with a variety of
condensed matter phenomena thou ; resta ; hatsugai ; bb since its inception
berry . Besides, various theoretical investigations, the geometric phases have
been experimentally tested in various cases, e.g. with photons p1 ; p2 ; p3 ,
with neutrons n1 ; n2 and with atoms a1 . The quantum state engineering of
cavity QED is in the state of art due to the rapid experimental/technological
progress of the subject. We hope that the theoretical scheme which we propose
for the laser field detuning induce quenching process in the dynamics of
geometric phase will be predicted by the experimental group based on photon.
The generation of the geometric phase (GP) is a witness of a singular point in
the energy spectrum that arises in all non-trivial geometric evolutions. In
this respect, the connection of the geometric phase with quantum phase
transition (QPT) has been explored very recently car ; zhu ; hamma ; sujgeo .
The geometric phase can be used as a tool to probe QPT in many body systems.
Since response times typically diverge in the vicinity of the critical point,
sweeping through the phase transition with a finite velocity leads to a
breakdown of adiabatic condition and generate interesting dynamical (non-
equilibrium ) effects. In the case of thermal phase transitions, the Kibble-
Zurek (KZ) mechanism kib ; zur explains the formation of defects via rapid
cooling. This idea of defect formation in second order phase transition has
been extended to zero temperature QPT zur1 ; dziar by studying the spin
models under linear quench. We will use this concept in the present study.
The micro-cavities of a photonic crystal are coupled through the exchange of
photons. Each cavity consists of one atom with three levels in the energy
spectrum, two of them are long lived and represent two spin states of the
system and the other represent excited states Ref. hart1 ; hart2 ; sujop ).
Externally applied laser and cavity modes couple to each atom of the cavity.
It may induce the Raman transition between these two long lived energy levels.
Under a suitable detuning between the laser and the cavity modes, virtual
photons are created in the cavity which mediate interactions with another atom
in a neighboring cavity. One can eliminate the excited states by choosing the
appropriate detuning between the applied laser and cavity modes. Then one can
achieve only two states per atom in the long lived state and the system can be
described by a spin-1/2 Hamiltonian hart1 ; hart2 ; sujop .
The Hamiltonian of our present study consists of three parts:
$H~{}=~{}{H_{A}}~{}+~{}{H_{C}}~{}+~{}{H_{AC}}$ (1)
Hamiltonians are the following
${H_{A}}~{}=~{}\sum_{j=1}^{N}{{\omega}_{e}}|e_{j}><e_{j}|~{}+~{}{\omega}_{ab}|b_{j}><b_{j}|$
(2)
where $j$ is the cavity index. ${\omega}_{ab}$ and ${\omega}_{e}$ are the
energies of the state $|b>$ and the excited state respectively. The energy
level of state $|a>$ is set as zero. $|a>$ and $|b>$ are the two stable state
of a atom in the cavity and $|e>$ is the excited state of that atom in the
same cavity. The following Hamiltonian describes the photons in the cavity,
${H_{C}}~{}=~{}{{\omega}_{C}}\sum_{j=1}^{N}{{c_{j}}}^{\dagger}{c_{j}}~{}+~{}{J_{C}}\sum_{j=1}^{N}({{c_{j}}}^{\dagger}{c_{j+1}}+h.c),$
(3)
where ${c_{j}}^{\dagger}({c_{j}})$ is the photon creation(annihilation)
operator for the photon field in the $j$’th cavity, ${\omega}_{C}$ is the
energy of photons and $J_{C}$ is the tunneling rate of photons between
neighboring cavities. The interaction between the atoms and photons and also
by the driving lasers are described by
${H_{AC}}~{}=~{}\sum_{j=1}^{N}[(\frac{{\Omega}_{a}}{2}e^{-i{{\omega}_{a}}t}+{g_{a}}{a_{j}})|e_{j}><a_{j}|+h.c]+[a\leftrightarrow
b].$ (4)
Here ${g_{a}}$ and ${g_{b}}$ are the couplings of the cavity mode for the
transition from the energy states $|a>$ and $|b>$ to the excited state.
${\Omega}_{a}$ and ${\Omega}_{b}$ are the Rabi frequencies of the lasers with
frequencies ${\omega}_{a}$ and ${\omega}_{b}$ respectively.
The authors of Ref. hart1 ; hart2 ; sujop have derived an effective spin
model by considering the following physical processes: A virtual process
regarding emission and absorption of photons between the two stable states of
neighboring cavity yields the resulting effective Hamiltonian as
${H_{xy}}=\sum_{j=1}^{N}B{{\sigma}_{j}}^{z}~{}+~{}\sum_{j=1}^{N}(\frac{J_{1}}{2}{{\sigma}_{j}}^{\dagger}{{\sigma}_{j+1}}^{-}~{}+~{}\frac{J_{2}}{2}{{\sigma}_{j}}^{-}{{\sigma}_{j+1}}^{-}+h.c)$
(5)
When $J_{2}$ is real then this Hamiltonian reduces to the XY model. Where
${{\sigma}_{j}}^{z}=|b_{j}><b_{j}|~{}-~{}|a_{j}><a_{j}|$,
${{\sigma}_{j}}^{+}=|b_{j}><a_{j}|$, ${{\sigma}_{j}}^{-}=|a_{j}><b_{j}|$ .
$\displaystyle H_{xy}$ $\displaystyle=$
$\displaystyle\sum_{i=1}^{N}(B{{\sigma}_{i}}^{z}~{}+~{}{J_{1}}({{\sigma}_{i}}^{x}{{\sigma}_{i+1}}^{x}+{{\sigma}_{i}}^{y}{{\sigma}_{i+1}}^{y})$
(6)
$\displaystyle+{J_{2}}({{\sigma}_{i}}^{x}{{\sigma}_{i+1}}^{x}-{{\sigma}_{i}}^{y}{{\sigma}_{i+1}}^{y}))$
$\displaystyle=\sum_{i=1}^{N}B({{\sigma}_{i}}^{z}~{}+~{}{J_{x}}{{\sigma}_{i}}^{x}{{\sigma}_{i+1}}^{x}~{}+~{}{J_{y}}{{\sigma}_{i}}^{y}{{\sigma}_{i+1}}^{y}).$
With ${J_{x}}=(J_{1}+J_{2})$ and ${J_{y}}=(J_{1}-J_{2})$.
We follow the references jame ; hart1 to present the analytical expression
for the different physical parameters of the system.
$B=\frac{\delta_{1}}{2}-\beta$ (7) $\displaystyle\beta$ $\displaystyle=$
$\displaystyle\frac{1}{2}[\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}({\Delta}_{b}-\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}-$
(8)
$\displaystyle\frac{{|{\Omega_{b}}|}^{2}}{4({\Delta}_{a}-{\Delta}_{b})}-{\gamma_{b}}{g_{b}}^{2}-{\gamma_{1}}{g_{a}}^{2}+{\gamma_{1}}^{2}\frac{{g_{a}}^{4}}{{\Delta_{b}}}-(a\leftrightarrow
b)]$
${J_{1}}=\frac{\gamma_{2}}{4}(\frac{{|{\Omega_{a}}|}^{2}{g_{b}}^{2}}{{{\Delta}_{a}}^{2}}+\frac{{|{\Omega_{b}}|}^{2}{g_{a}}^{2}}{{{\Delta}_{b}}^{2}})$
(9)
${J_{2}}=\frac{\gamma_{2}}{2}(\frac{{\Omega_{a}}{\Omega_{b}}g_{a}g_{b}}{{\Delta}_{a}{\Delta_{b}}}).$
(10)
Where $\gamma_{a,b}=\frac{1}{N}\sum_{k}\frac{1}{{\omega}_{a,b}-{\omega}_{k}}$
$\gamma_{1}=\frac{1}{N}\sum_{k}\frac{1}{({\omega}_{a}+{\omega}_{b})/2-{\omega}_{k}}$
and
$\gamma_{2}=\frac{1}{N}\sum_{k}\frac{e^{ik}}{({\omega}_{a}+{\omega}_{b})/2-{\omega}_{k}}$
${\delta_{1}}={\omega}_{ab}-({\omega}_{a}-{\omega}_{b})/2$,
${\Delta}_{a}={\omega}_{e}-{\omega}_{a}$.
${\Delta}_{b}={\omega}_{e}-{\omega}_{a}-({\omega}_{ab}-{\delta_{1}})$.
${{\delta}_{a}}^{k}={\omega}_{e}-{\omega}_{k}$,
${{\delta}_{b}}^{k}={\omega}_{e}-{\omega}_{k}-({\omega}_{ab}-{\delta_{1}})$,
$g_{a}$ and $g_{b}$ are the couplings of respective transition to the cavity
mode, ${\Omega}_{a}$ and ${\Omega}_{b}$ are the Rabi frequency of laser with
frequency $\omega_{a}$ and $\omega_{b}$.
Model Hamiltonian and Quantum Phases: We express our model Hamiltonian in a
more explicit way as
$H~{}=~{}\sum_{n}~{}[(1+\alpha)~{}S_{n}^{x}S_{n+1}^{x}~{}+~{}(1-\alpha)~{}S_{n}^{y}S_{n+1}^{y}+B\sum_{n}~{}S_{n}^{z}$
(11)
where $S_{n}^{\alpha}$ are the spin-1/2 operators. We assume that the $XY$
anisotropy $0<\alpha\leq 1$ and magnetic field strength is $h\geq 0$. The
parameters correspondence between the micro cavities and spin chain are the
following: ${J_{1}=1}$ and ${J_{2}}=\alpha$.
Figure 1: Color online, variation of geometric phase with time for different
quenching time for $k=\pi/3$. The anisotropy parameter ($\alpha=1$, transverse
Ising model) for the dash curves and inset for $\alpha=0$ (isotropic XX spin
model). Figure 2: Color online, variation of the geometric phase with time for
three different quenching time for $k=\pi/3$. The solid curve is for the
$\alpha=0.2$ and dashed curve is for $\alpha=0.8$. Figure 3: Color online,
variation of geometric phase with time for the quenching time
${\tau_{q}=1000}$ for different anisotropy in exchange couplings for
$k=\pi/3$. The all curves are meeting in the point at $t={\tau_{q}}$.
Geometric Phase and Criticality: Here, we calculate the geometric phase and
its dynamics under the quenching of magnetic field. In this model, the
geometric phase of the ground state is evaluated by applying a rotation of
$\phi$ around the $z$-axis in a closed circuit to each spin car ; car1 . A new
set of Hamiltonians $H_{\phi}$ is constructed from the Hamiltonian (H) as
$H_{\phi}=U(\phi)~{}H~{}U^{\dagger}(\phi)$ (12)
where $U(\phi)=\prod_{j=-M}^{+M}\exp(i\phi\sigma_{j}^{z}/2)$ and
$\sigma_{j}^{z}$ is the $z$ component of the standard Pauli matrix at site
$j$. Here $M$ is the integer which relates with the lattice site numbers by
the following relation $2M+1=N$. In our atom-cavity system, the rotation of
laser field around the z-axis is equivalent to the rotation of the quantum
spin system around the z-axis. The family of Hamiltonians generated by varying
$\phi$ has the same energy spectrum as the initial Hamiltonian and $H(\phi)$
is $\pi$-periodic in $\phi$. The ground state $|g>$ of the system is expressed
as
$|g>=\prod_{k>0}(\cos\frac{\theta_{k}}{2}|0>_{k}|0>_{-k}-i\exp(2i\phi)\sin\frac{\theta_{k}}{2}|1>_{k}|1>_{-k})$
(13)
where $|0>_{k}$ and $|1>_{k}$ are the vacuum and single fermionic excitation
of the $k$-th momentum mode respectively. The angle $\theta_{k}$ is given by
$\cos\theta_{k}=\frac{\cos k-B}{\Lambda_{k}}$ (14)
and $\displaystyle{\Lambda_{k}=\sqrt{(\cos k-B)^{2}+\alpha^{2}\sin^{2}k}}$ is
the energy gap above the ground state. The ground state is a direct product of
$N$ spins, each lying in the two-dimensional Hilbert space spanned by
$|0>_{k}|0>_{-k}$ and $|1>_{k}|1>_{-k}$. For each value of $k$, the state in
each of the two dimensional Hilbert space can be represented as a Bloch vector
with coordinates $(2\phi,\theta_{k})$. In our previous study sujop ; bonodi
we have shown explicitly the derivation of geometric phase for this problem.
$\Gamma_{k}=\pi(1-\cos\theta_{k})$ (15)
The total geometric phase of the state $|g>$ is given by
$\Gamma_{g}=\sum_{k}~{}\Gamma_{k}$ . For an adiabatic evolution, if the
initial state is an eigenstate, the evolved state remains in the eigenstate.
Now we derive the instantaneous geometric phases of this system due to a
gradually decreasing magnetic field, i.e., the quenching field.
$B(t<0)=-\frac{t}{\tau_{q}}$ (16)
$B(t)$, driving the transition, is assumed to be linear with an adjustable
time parameter $\tau_{q}$ ($1/{\tau_{q}}$ is the quenching rate). The system
be initially at time $t(<0)<<\tau_{q}$ such that $B(t)>>1$, in our parameter
space of cavity arrays system,
${\delta_{1}}(t)>>[\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}({\Delta}_{b}-\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}-\frac{{|{\Omega_{b}}|}^{2}}{4({\Delta}_{a}-{\Delta}_{b})}-{\gamma_{b}}{g_{b}}^{2}-{\gamma_{1}}{g_{a}}^{2}+{\gamma_{1}}^{2}\frac{{g_{a}}^{4}}{{\Delta_{b}}}-(a\leftrightarrow
b)]$. Here we consider the quenching of the magnetic field by considering
quenching the detuning field ${\delta}_{1}$.
The instantaneous ground state at any instant $t$ is given by
$|\psi_{0}(t)>=\prod_{k}(\cos\frac{\theta_{k}(t)}{2}|0>_{k}|0>_{-k}-i\exp(2i\phi)\sin\frac{\theta_{k}(t)}{2}|1>_{k}|1>_{-k})$
(17)
We now use eqn. (14) and (13) to derive the geometric phase of the $k^{th}$
mode which yields
$\Gamma_{k}(t)=\pi\left(1-\frac{\cos k+\frac{t}{\tau_{q}}}{\sqrt{(\cos
k+\frac{t}{\tau_{q}})^{2}+\alpha^{2}\sin^{2}k}}\right)$ (18)
The geometric phase for an isotropic system with $\alpha=0$ and transverse
Ising model with $\alpha=1$ may now be easily obtained.
$\displaystyle\rm{For}~{}\alpha=0,~{}~{}~{}~{}\Gamma_{k}(t)$ $\displaystyle=$
$\displaystyle 2\pi\Theta(|t|-{{\tau}_{q}})$
$\displaystyle\rm{and}~{}\rm{for}~{}\alpha=1,~{}~{}\Gamma_{k}(t)$
$\displaystyle=$ $\displaystyle\pi\left(1-\frac{\cos
k+\frac{t}{\tau_{q}}}{\sqrt{1+\frac{t^{2}}{\tau_{q}^{2}}+2\frac{t}{\tau_{q}}\cos
k}}\right)$ (19)
For ${\alpha}=0$, one of the Rabi frequencies, either (${\Omega_{a}}$) or
(${\Omega_{b}}$) is zero. For ${\alpha}=1$, where both of Rabi frequencies are
non zero. Finally, at $t=0$, when the magnetic field is gradually turned off,
the situation is a bit different. The configuration of the final state will
depend on the number of kinks generated in the system due to phase transition
at or near $t=-{\tau_{q}}$ and as such it will depend on the quench time
$\tau_{q}$ zur . The number of kinks is the number of quasi-particles excited
at $B=0$ ( in our parameter space of micro-cavities, ${{\delta}_{1}}/2=\beta$,
and the quasiparticle excitation is the excitation of polariton) and is given
by ${\cal{N}}={\sum_{k}}p_{k}$ where $p_{k}$ is the excitation probability
(for the slow transition) and is given by the Landau Zener formula zener
$p_{k}\approx\exp{(-2\pi\tau_{q}k^{2})}$. As different pairs of quasi-
particles ($k,-k$) evolve independently for large values of $\tau_{q}$, it is
likely that only one pair of quasi-particles with momenta $(k_{0},-k_{0})$
will be excited. Where $k_{0}(=\frac{\pi}{N})$ corresponds to the minimum
value of the energy $\Lambda_{k}$. Thus the condition for adiabatic transition
in a finite chain is given by $\tau_{q}>>\frac{N^{2}}{2\pi^{3}}$. Hence, well
in the adiabatic regime, the final state at $t=0$ is given by
$\displaystyle|\psi_{final}>$ $\displaystyle=$
$\displaystyle|1>_{k_{0}}|0>_{-k_{0}}$ (20) $\displaystyle\prod_{k,k\neq\pm
k_{0}}(\cos\frac{\theta_{k}}{2}|0>_{k}|0>_{-k}$
$\displaystyle-i\exp(2i\phi)\sin\frac{\theta_{k}}{2}|1>_{k}|1>_{-k})$
Figure 4: Color online, variation of the geometric phase with time for two
different quenching time, $\tau_{q}=100,10,000$ for four different anisotropy
parameters $\alpha=0.2,0.5,0.8,1.0$ for the blue, cyan, red, green
respectively. Here $k=\pi/3$. Figure 5: Color online, this three dimensional
figures shows the variation of first order derivative of geometric phase with
time, with the anisotropy exchange parameter and time. Its shows the existence
of XX criticality (${\alpha=0}$ ).
This state is similar to the direct product of only $N-1$ spins oriented along
$(2\phi,\theta_{k})$ where the state of the spin corresponding to momentum
$k_{0}$ does not contribute to the geometric phase. The total geometric phase
of this state is given by $\Gamma_{final}(t=0)=\sum_{k,k\neq\pm
k_{0}}\pi(1-\cos\theta_{k})$. For isotropic, $\alpha=0$ and for quantum spin
chain system, geometric phase shows a sharp drop and finally become zero at
$t={\tau_{q}}$ (see the inset of the Fig. 1) for this case optical cavity
array system has only single Rabi frequency oscillation. In our study, we set
the energy scale $J_{1}=1$, that leads to the following relations in the Rabi-
frequencies, laser frequencies and coupling strengths in the following way,
$\frac{4}{\gamma_{2}}=(\frac{{|{\Omega_{a}}|}^{2}{g_{b}}^{2}}{{{\Delta}_{a}}^{2}}+\frac{{|{\Omega_{b}}|}^{2}{g_{a}}^{2}}{{{\Delta}_{b}}^{2}})$.
We achieve the isotropic $H_{XX}$ Hamiltonian, either ${\Omega}_{a}$ or
${\Omega}_{b}$ is zero. For ${\Omega}_{a}=0$, the reduce analytical expression
for the magnetic field and $J_{1}$ exchange are
$B=\frac{\delta_{1}}{2}-\frac{1}{2}[\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}({\Delta}_{b}-\frac{{|{\Omega_{b}}|}^{2}}{4{\Delta}_{b}}-\frac{{|{\Omega_{b}}|}^{2}}{4({\Delta}_{a}-{\Delta}_{b})}-{\gamma_{b}}{g_{b}}^{2}-{\gamma_{1}}{g_{a}}^{2}+{\gamma_{1}}^{2}\frac{{g_{a}}^{4}}{{\Delta_{b}}}]$.
and
${J_{1}}=\frac{\gamma_{2}}{4}(\frac{{|{\Omega_{b}}|}^{2}{g_{a}}^{2}}{{{\Delta}_{b}}^{2}})$
respectively. For ${\Omega_{b}}=0$, the analytical expression for $B$ and
$J_{1}$ can be obtained from Eq. (7) and Eq. (8) by putting ${\Omega_{b}}=0$.
For $\alpha=1$, the spin chain is in a transverse magnetic field. One can
achieve this limit by adjusting the detuning field, Rabi frequencies and the
coupling strengths. For this limit, the analytical expression between the Rabi
frequency and coupling strength is
${\Delta}_{b}=\frac{\gamma_{2}\Omega_{a}\Omega_{b}g_{a}g_{b}}{2\Delta_{a}}$,
two Rabi frequency oscialltion exist for this case. The analytical expression
for effective magnetic filed and the $J_{1}$ coupling can be obtained by
substituting the analytical expression of $\Delta_{b}$ in Eq (7) and Eq. (8).
The dynamics of the geometric phase shows no sharp drop of the geometric phase
through out the quenching period. The dynamics of the geometric phase under
quenching is different for isotropic quantum spin chain system ($H_{XX}$).
Inset of the Fig. 1 shows the dynamics of geometric phase for $H_{XX}$ model.
We will also observe that the first order derivative of the geometric phase
also behave differently for these two limit of our model Hamiltonian.
Therefore it is very clear from this study that the presence of single Rabi
frequency and two Rabi frequency oscillation in the system of cavity QED
changes the dynamics of geometric phase differently.
In Fig. 2, we present the study of the geometric phase for different values of
anisotropy exchange interactions ${\alpha}=0.2,0.8$ and for three different
quenching times ${\tau}_{q}=250,500,1000$, which can be achieved by
manipulating two Rabi frequencies ${\Omega}_{a}$ and ${\Omega}_{b}$. It
reveals from our study that the qualitative behavior of the geometric phase is
same but for the lower value of exchange anisotropy decaying rate is higher
than for the higher values of exchange anisotropy. The geometric phase for
different anisotropic exchange parameters crosses at the point $t={\tau_{q}}$
for the same ${\tau}_{q}$. This common crossing point is nothing but the
quantum critical point for the system at $B=1$ for the same $\tau_{q}$. This
has shown explicitly in Fig. 5 of this study.
Apart from that the geometric phase for lower value of $\alpha$ also crosses
the geometric phases of other $\tau_{q}$ but at the time lower than
$\tau_{q}$. Therefore it is clear from this study due to the smaller values of
Rabi frequencies (${\Omega}_{a}$ or ${\Omega_{b}}$) decaying rate of geometric
phase is sharper compare to the larger values of Rabi frequencies. The same
physics can also be obtained by regulating ${\Delta}_{a}$ and ${\Delta}_{b}$.
It is very clear from the analytical expression of Eq. 14 and Eq. 15 that the
dynamics of geometric phase depends on the magnetic field (time dependent
quenching field) and the exchange anisotropy, therefore for a fixed $\tau_{q}$
the common meeting point of the geometric phase occurs at the quenching time.
In Fig.3, we study the effect of different exchange interactions for
$\tau_{q}=1000$. We observe that the all geometric phases crosses at the point
$t={\tau}_{q}$. It reveals from this study that for a fixed quenching rate,
the different exchange anisotropy (the presence of different Rabi frequencies
oscillation) has no effect to change the common crossing point of the
geometric phase.
In our study, $1/{\tau_{q}}$, is the quenching rate for smaller values of
$\tau_{q}$ is larger and smaller for higher values of $\tau_{q}$. In Fig. 4
shows the dynamical evalution of geometric phase for higher and lower
quenching rate. It appears from our study that for higher quenching rate
$(\tau_{q}=100)$ all curves for different anisotropy. For lower quenching rate
is smaller for this case there is no collapse of the dynamical behaviour of
the geometric phase.
In Fig. 4, we study the effect of slower and rapid quenching rate on the
dynamical behaviour of geometric phase. It appears from our study that for
higher quenching rate ($\tau_{q}=100$) all curves for different anisotropy
merge into a single line as if the exchange anisotropy,i.e., the different
values of two Rabi frequency oscillations has no effect on the dynamical
behaviour of the geometric phase. The dynamical behaviour of the geometric
phase for lower quenching rate ($\tau_{q}=10,000$) is the same as that of
$\tau_{q}=1000$. The crossing point of the geometric phase at $B=1$ for a
fixed $\tau_{q}$ is the same for all quenching rate.
To study the criticality of the geometric phase, we study the variation of the
geometric phase (${\Gamma}_{k}$) and its derivative with respect to the quench
induced magnetic filed ($B$) i.e. $({\frac{d{\Gamma}_{k}}{dB}})$ with time. We
find the non-analytic behavior of the derivative at $t={\tau}_{q}$. The
analytical expression for the derivative is
${\frac{d{\Gamma}_{k}(t)}{dB}}=\frac{\pi{\alpha}^{2}sin^{2}(k)}{{({({cos(k)+t/{{\tau}_{q}})}^{2}+{\alpha}^{2}sin^{2}(k)})}^{3/2}}$
(21)
Fig. 5, shows the total variation of ${\frac{d{\Gamma}_{k}(t)}{dB}}$. The non-
analytical behavior for $\alpha=0$ at $t=-{\tau_{q}}$ helps us to predict XX
criticality under the linear quenching process. The analysis with different
values of $\tau_{q}$ shows that the appearance of XX criticality is
independent of $\tau_{q}$, i.e., independent of fast and slow quenching rates.
Therefore, it is clear from our study that the quantum criticality of the
geometric phase only obtain in the presence of single Rabi frequency
oscillation in the system.
Conclusions: We have studied the quantum criticality of the geometric phase
for optical cavity arrays. We have predicted XX criticality of the geometric
phase, which is independent of the quenching rate. We present the result for
the dynamics of geometric phase for different quenching rate. The presence of
two Rabi frequencies oscillation wash out the quantum critical behavior of the
geometric phase.
Acknowledgments
The author would like to thank, Prof. S. Girvin for useful discussions during
the international workshop/school on Dirac Materials at ICTS (December, 2012)
and also the library of Raman Research Institute (Mr. Manjunath). Finally, the
author would like to thank Dr. P. K. Mukherjee for reading the manuscript
carefully.
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|
arxiv-papers
| 2013-09-30T09:29:37 |
2024-09-04T02:49:51.749809
|
{
"license": "Public Domain",
"authors": "Sujit Sarkar",
"submitter": "Sujit Sarkar",
"url": "https://arxiv.org/abs/1310.0003"
}
|
1310.0056
|
# modernizing PHCpack through phcpy
Jan Verschelde∗† * Corresponding author: [email protected]† University of
Illinois at ChicagoCopyright © 2014 Jan Verschelde. This is an open-access
article distributed under the terms of the Creative Commons Attribution
License, which permits unrestricted use, distribution, and reproduction in any
medium, provided the original author and source are credited.
http://creativecommons.org/licenses/by/3.0/
###### Abstract
PHCpack is a large software package for solving systems of polynomial
equations. The executable phc is menu driven and file oriented. This paper
describes the development of phcpy, a Python interface to PHCpack. Instead of
navigating through menus, users of phcpy solve systems in the Python shell or
via scripts. Persistent objects replace intermediate files.
## 1 Introduction
Our mathematical problem is to solve a system of polynomial equations in
several variables. The discrete part of the output data includes the number of
solutions and degrees of positive dimensional solution sets. When the input is
exact or if the coefficients of the polynomial can be given with any accuracy,
the isolated solutions can be then approximated up to any accuracy and for
each positive dimensional solution component, as many generic points as its
degree can be computed.
Version 1.0 of PHCpack was archived in [Ver99]. PHCpack incorporates two
external software packages: MixedVol [GLW05] and QDlib [HLB01]. Although the
original focus was to approximate all isolated complex solutions, PHCpack
prototyped many of the early algorithms in numerical algebraic geometry
[SVW03], [SVW05]. Recent updates are listed in [Ver10].
The Python interface to PHCpack got a first start when Kathy Piret met William
Stein at the software for algebraic geometry workshop at the Institute for
Mathematics and its Applications in Minneapolis in the Fall of 2006. The first
version of this interface is described in [Pir08]. Sage [S+13] offers the
interface phc.py, developed by William Stein, Marshall Hampton and Alex
Jokela. Version 0.0.1 of phcpy originated at lecture 40 of the author in the
graduate course MCS 507 in the Fall of 2012, as an illustration of Sphinx
[Bra]. Version 0.1.0 was prepared for presentation at EuroSciPy 2013 (August
2013), version 0.1.4 corresponds to the first version of this paper. The
current version of phcpy is 0.1.5.
We first outline in the next section the application of numerical homotopy
continuation methods to compute, over the complex numbers, all isolated
solutions and all positive dimensional irreducible solution sets of a
polynomial system. Then we describe how phcpy relates to other interfaces to
PHCpack. The functionality of phcpy is then summarized briefly as the online
Sphinx documentation is more extensive and still growing. The Python interface
to PHCpack builds directly on the C interface to the Ada code.
Related software packages that apply homotopy continuation methods to solve
polynomial systems are (in alphabetical order): Bertini [BHSW], [BHSW08],
HOMPACK90 [WSM+97] (the successor of HOMPACK [WBM87]), HOM4PS [GLL02],
HOM4PS-2.0 [LLT08], NAG4M2 [Ley11], PHoM [GKK+04], and pss3.0.5 [Mal]. As
polynomial homotopy continuation methods involve many algorithms from various
fields of computing, every software has its unique strengths and the statement
“ _no one package provides all such capabilities_ ” quoted from [BHSW08]
remains true today.
## 2 Polynomial Homotopy Continuation
Our mathematical problem is to solve a polynomial system $f({\bf x})={\bf 0}$
in several variables ${\bf x}=(x_{1},x_{2},\ldots,x_{n})$. A homotopy connects
$f({\bf x})={\bf 0}$ to a system $g({\bf x})={\bf 0}$ with known solutions:
$h({\bf x},t)=\gamma(1-t)g({\bf x})+tf({\bf x})={\bf
0},\quad\gamma\in{\mathbb{C}}.$
For almost all values for $\gamma$, the solutions of $h({\bf x},t)={\bf 0}$
are regular for all $t\in[0,1)$. Numerical continuation methods track solution
paths defined by $h({\bf x},t)={\bf 0}$.
For systems with natural parameters $\lambda$, we solve
$f({\mbox{\boldmath$\lambda$}},{\bf x})={\bf 0}$ first for generic values of
the parameters ${\mbox{\boldmath$\lambda$}}={\mbox{\boldmath$\lambda$}}_{0}$
and then use
$h({\bf x},t)=(1-t)f({\mbox{\boldmath$\lambda$}}_{0},{\bf
x})+tf({\mbox{\boldmath$\lambda$}}_{1},{\bf x})={\bf 0},$
to solve a specific instance $f({\mbox{\boldmath$\lambda$}}_{1},{\bf x})={\bf
0}$.
The schematic in Figure 1 illustrates that singular solutions along the paths
are avoided by a generic choice of the parameters
${\mbox{\boldmath$\lambda$}}_{0}$ at $t=0$.
Figure 1: A generic choice for start parameters avoids singularities along the
paths.
Solving a system that has solution sets of positive dimension, e.g.: a curve
or surface, means to compute its dimension and its degree. For a solution set
of dimension $d$, we add to the system $d$ linear equations with random
coefficients to reduce the problem to computing isolated solutions. If the
dimension is not known in advance, then it can be computed in a cascade of
homotopies [SV00]. Because the coefficients of the linear equation are random,
the solutions of the system that satisfy the random linear equations will be
isolated. If the solution set has multiplicity one, the solutions of the
augmented system will be isolated points of multiplicity one. Moreover, the
number of isolated solutions of the augmented system that lie on the
$d$-dimensional solution set of the original system will be equal to the
degree of the $d$-dimensional solution set. Thus a positive dimensional
solution set of dimension $d$ is represented by a set of $d$ random linear
equations and as many points in the intersection of the original system with
those random linear equations as the degree of the $d$-dimensional solution
set. In numerical algebraic geometry, this representation is called a witness
set.
For sparse polynomial systems with very few monomials appearing with nonzero
coefficient (in an extreme case, we consider binomial systems that have
exactly two monomials with nonzero coefficient in each equation), we can
represent positive dimensional solution sets by monomial maps. For example,
the two equations $x^{2}y-zx=0$, $x^{2}z-y^{2}x=0$ have as solutions three
monomial maps: $(x=0,y=\lambda_{1},z=\lambda_{2})$,
$(x=\lambda_{1},y=\lambda_{1}^{2},z=\lambda_{1}^{3})$, and
$(x=\lambda_{1},y=0,z=0)$, for parameters $\lambda_{1}$ and $\lambda_{2}$.
These monomial maps form the leading terms of Puiseux series developments for
general algebraic sets.
Surveys on homotopy continuation are [AG93], [AG97], [Li03], and [Wat86],
[Wat89], [Wat02]. Book treatments are in [AG03], [Mor87], and [SW05].
## 3 Interfaces to PHCpack and phc
This paper is mainly concerned with software problems. There are at least
three motivations to develop phcpy:
* 1.
PHCpack is a large Ada package, its executable phc operates via menus, with
input and output to files. With phcpy we provide an interpreter interface to
phc.
* 2.
The code in PHCpack lacks adequate _user_ documentation so that many of its
features are not obviously accessible to users. The Python modules of phcpy
refactor the functionality of PHCpack and beautiful documentation is generated
by Sphinx [Bra].
* 3.
As many new algorithms were first implemented with PHCpack, reproducibility
[SBB13] of published computational results can be automated via regression
tests with Python scripts.
Because also other interfaces to PHCpack may accomplish the same goals
outlined above, we first give an overview of the interfaces to PHCpack.
The first interface to PHCpack was based on the OpenXM [MNO+11] protocol for
the interaction of software components. The virtue of this protocol is that
only an executable version of the software is required and one does not need
to compile the code.
The interfaces to PHCpack from Maple [LV04], MATLAB & Octave [GV08b], and
Macaulay2 [GPV13] only require the executable phc. This type of interface
works in three stages: (1) prepare an input file for phc; (2) call phc with
some options, the input file, and the name of an output file; (3) parse the
output file to extract the results. In principle, everything that can be done
via the command-line menus of phc can thus also be performed via Maple
procedures, MATLAB, Octave, or Macaulay2 scripts.
Figure 2 shows the interfaces to PHCpack.
Figure 2: Diagram of the interfaces to PHCpack and phc. The interfaces
PHCpack.m2, PHCmaple, PHClab, depicted to the right of the antidiagonal line
require only the executable version phc. The other interfaces PHClib, PHCmpi,
and phcpy are based on the source code of PHCpack.
The C interface to PHCpack, described in [LV06], offers the C programmer
access to the path trackers of PHCpack. This interface was developed for use
with the Message Passing Interface [SOH+98] and serves also as the basis for
phcpy. In the C interface, the data structures for polynomials and solutions
are not duplicated. Instead of data structure duplication, one can enter into
the C interface routine a polynomial term after term. The interface then
behaves like a state machine.
Why would phcpy be any better than the other interfaces? Leaving aside the
growing popularity of Python for scientific computing, the replacement of
files by persistent objects enabled the implementation of a generator for the
path trackers. After initialization of the homotopy (with target, start
system, and one start solution), the user can call a “next” function to
compute the next point at the solution path that originates at the start
solution given at initialization. This “next” function (available for standard
double, double double, quad double precision, and arbitrary multiprecision)
allows a detailed investigation of the properties of a particular solution
path. In addition, it gives the user a fine control over the order of
execution. If desired, the tolerances and the step size can be adjusted as
needed in an application that plots solution trajectories.
Another (future) application of phcpy is a web interface, such as at
https://kepler.math.uic.edu (beta version) presented by Xiangcheng Yu at the
SIAM AG 2013 conference in the first week of August 2013.
## 4 Using phcpy
The blackbox solver of PHCpack is its most widely used function. In phcpy,
this blackbox solver is available in the function solve of the module solver.
The solver takes on input a list of strings that contain valid representations
of polynomials. On return is a list of strings, which contain the solutions of
the system.
>>> from phcpy.solver import solve>>> from phcpy.phcpy2c import
py2c_set_seed>>> f = ["x**2*y**2 + x + y;","x*y + x + y + 1;"]>>>
py2c_set_seed(21320)0>>> s = solve(f,silent=True)>>> len(s)4>>> print s[0]t :
1.00000000000000E+00 0.00000000000000E+00m : 1the solution for t :x :
-1.00000000000000E+00 0.00000000000000E+00y : -1.61803398874989E+00
0.00000000000000E+00== err : 2.143E-101 = rco : 4.775E-02 = res : 2.220E-16
=With py2c_set_seed() we fix the seed of the random number generator for the
coefficients of the start system in the homotopy, which makes for predictable
runs. Otherwise, the solve() each time generates different coefficients in the
homotopies and the order of the solutions on return may differ. For each
solution, the triplet (err,rco,res) indicates the quality of the solution:
* •
err: the norm of the last update made by Newton’s method (forward error),
* •
rco: estimate for the inverse condition number of the Jacobian matrix,
* •
res: norm of the evaluated solution (backward error).
With double double and quad double arithmetic we get more accurate solutions.
To predict the number of isolated solutions with the mixed volume:
>>> from phcpy.solver import mixed_volume>>> mixed_volume(f)4Version 0.1.5 of
phcpy contains the following modules:
* •
solver: a blackbox solver, mixed-volume calculator, linear-product root count
and start system, path trackers, deflation for isolated singular solutions.
* •
examples: a selection of interesting benchmark systems. Typing python
examples.py at the command prompt calls the blackbox solver on all benchmark
examples, thus providing an automatic regression test.
* •
families: some problems can be formulated for any number of variables.
* •
phcmaps: monomial maps as solutions of binomial systems.
* •
phcsols: conversion of PHCpack solution strings into Python dictionaries.
* •
phcsets: basic tools to manipulate positive dimensional solution sets.
* •
phcwulf: basic client/server setup to solve many systems.
* •
schubert: the Pieri homotopies solve particular polynomial systems arising in
enumerative geometry.
The number of exported functions, documented by Sphinx [Bra] runs in the
several hundreds. The code of version 0.1.1 of phcpy was improved with the aid
of Pylint [The], yielding a global rating of 9.73/10.
## 5 The Design of phcpy
The design of phcpy is drawn in Figure 3. This design can be viewed as an
application of a façade pattern (see Figure B.31 in [Bai08]). The façade
pattern plays a strong role in converting legacy systems incrementally to more
modern software and is appropriate as phcpy should be viewed as a
modernization of PHCpack. The implementation of use_c2phc.adb applies the
chain of responsibility pattern (see Figure B.12 in [Bai08]), calling handlers
to specific packages in PHCpack. That we use the name phcpy and not PyPHC
indicates that phcpy is more than just an interface.
Figure 3: The design of phcpy depends on PHClib, a library of various
collections of C functions, through one file phcpy2c.c (with documentation in
the corresponding header phcpy2c.h) which encodes the Python bindings. PHClib
interfaces to the Ada routines of PHCpack through one Ada procedure
use_c2phc.adb. The collection of parallel programs (MPI2phc) using message
passing (MPI) depends on PHClib.
The code for phcpy builds directly on the C interface to PHCpack. The C
interface was developed to use the Message Passing Interface (MPI) [SOH+98].
In joint work with Yusong Wang [VW02], [VW04a], [VW04b], Yan Zhuang [VZ06],
Yun Guan [GV08a], and Anton Leykin [LV05], [LV09], [LVZ06], the main program
was always a C program. The C interface described in [LV06] is centered around
one gateway function use_c2phc. To the Ada programmer, this function has the
specification
function use_c2phc ( job : integer; a : C_intarrs.Pointer; b :
C_intarrs.Pointer; c : C_dblarrs.Pointer ) return integer;The prototype of the
corresponding C function is
extern int _ada_use_c2phc ( int task, int *a, int *b, double *c );With
use_c2phc we obtain one uniform streamlined design of the interface: the C
programmer calls one single Ada function _ada_use_c2phc. What use_c2phc
executes depends on the job number. The (a,b,c) parameters are flexible enough
to pass strings and still provide some form of type checking (which would not
be possible had we wiped out all types with void*).
To make _ada_use_c2phc usable, we have written a number of C wrappers,
responsible for parsing the arguments of the C functions to be passed to
_ada_use_c2phc. The extension module and the shared object for the
implementation of phcpy is a set of wrappers defined by phcpy2c.c and
documented by phcpy2c.h. As a deliberate design decision of phcpy, all calls
to functions in PHCpack pass through the C interface. By this design, the
development of phcpy benefits the C and C++ programmers.
## 6 Obtaining, Installing, and Contributing
PHCpack and phcpy are distributed under the GNU GPL license (version 2 or any
later version). Recently a new repository PHCpack was added on github with the
source code of version 2.3.85 of PHCpack, which contains version 0.1.5 of
phcpy. Executable versions for Linux, Mac, and Windows are available via the
homepage of the author.
The code was developed on a Red Hat Enterprise Linux Workstation (Release 6.4)
and a MacBook Pro laptop (Mac OS X 10.8.5) using the GNAT GPL 2013 compiler.
Versions 2.6.6 and 2.7.3 of Python, respectively on Linux and Mac, were used
to develop phcpy. Packaged binary distributions of phcpy for the platforms
listed above are available via the homepage of the author.
Although the blackbox solver of PHCpack has been in use since 1996, phcpy
itself is still very much in beta stage. Suggestions for improvement and
contributions to phcpy will be greatly appreciated.
## 7 Acknowledgments
The author thanks Max Demenkov for his comments and questions at the poster
session at EuroSciPy 2013. In particular the question on obtaining all
solutions along a path led to the introduction of generator functions for the
path trackers in version 0.1.4 of phcpy.
This material is based upon work supported by the National Science Foundation
under Grant No. 1115777.
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|
arxiv-papers
| 2013-09-30T21:05:03 |
2024-09-04T02:49:51.759015
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Jan Verschelde",
"submitter": "Jan Verschelde",
"url": "https://arxiv.org/abs/1310.0056"
}
|
1310.0104
|
# Pure Subspaces, Generalizing the Concept
of Pure Spinors
Carlos Batista
Departamento de Física
Universidade Federal de Pernambuco
50670-901 Recife-PE, Brazil
[email protected]
###### Abstract
The concept of pure spinor is generalized, giving rise to the notion of pure
subspaces, spinorial subspaces associated to isotropic vector subspaces of
non-maximal dimension. Several algebraic identities concerning the pure
subspaces are proved here, as well as some differential results. Furthermore,
the freedom in the choice of a spinorial connection is exploited in order to
relate twistor equation to the integrability of maximally isotropic
distributions. (Keywords: Isotropic Spaces, Pure Spinors, Integrability,
Clifford Algebra, Twistors)
## 1 Introduction
It is well-known that given a spinor $\hat{\varphi}$ one can construct a
vector subspace $N_{\hat{\varphi}}$ spanned by the vectors that annihilate
$\hat{\varphi}$ under the Clifford action, $v\in N_{\hat{\varphi}}$ whenever
$v\cdot\hat{\varphi}=0$. This subspace is necessarily isotropic, which means
that $\langle v,v\rangle=0$ for all $v\in N_{\hat{\varphi}}$. Particularly, if
the dimension of $N_{\hat{\varphi}}$ is maximal $\hat{\varphi}$ is said to be
a pure spinor. Thus, to every spinor it is associated an isotropic vector
subspace. The idea of the present article is to go the other way around and
associate to every isotropic subspace $I$ a spinorial subspace $\hat{L}_{I}$.
As we shall see, these spinorial subspaces provide a natural generalization
for the concept of pure spinors and, therefore, we shall say that
$\hat{L}_{I}$ is the pure subspace associated to $I$. As a consequence, some
classical results are shown to be just particular cases of broader theorems.
Hence, this study can shed more light on the role played by the pure spinors
on physics and mathematics.
Spinors have been around for a long time, they were first discovered a century
ago by Élie Cartan [1]. Since then, they led to great achievements in a
multitude of areas of mathematics and physics. For instance, spinors are of
fundamental importance in particle physics, since fermions are represented by
spinorial fields (see [2] for an alternative treatment). In addition, spinors
can often be used as helpful tools for calculating physical and geometrical
quantities. As examples let us recall Witten’s elegant proof of energy
positivity in general relativity [3] and the calculation of scattering
amplitudes in quantum field theory [4]. From the geometrical point of view, a
particularly important class of spinors is formed by the so-called pure
spinors. Such spinors are also acquiring increasing significance in high-
energy physics, as illustrates the pure spinor formalism in string theory [5].
In spite of such undeniable usefulness, the role played by spinors is
certainly not completely understood yet. Hopefully, the work presented here
will contribute for a further understanding of spinorial calculus and,
specially, of the meaning of pure spinors.
The isotropic subspaces are of great relevance in mathematical-physics, a fact
that became clear after the works by R. Penrose wherein he introduced the
spinorial calculus and the null tetrad formalism in 4-dimensional general
relativity [6]. Such approaches, in which isotropic distributions play a
prominent role, brought much progress to general relativity. Notably,
Kinnersley was able to find analytically all solutions of Einstein’s vacuum
equation in 4 dimensions for the case of the space-time admitting two
independent integrable distributions of isotropic planes [7, 8, 9].
Particularly, all known 4-dimensional black-holes are contained in this class
of solutions. Higher-dimensional manifolds have been a subject of great
importance in theoretical physics, not to mention in mathematics. Following
the successful track of 4-dimensional general relativity recent works have
shown that isotropic structures do also play an important role in higher
dimensions. In reference [10] it was proved a relation between integrable
isotropic structures and the existence of a conformal Killing-Yano tensor,
while in [11] these structures were used to classify the Weyl tensor as well
as to partially generalize the Goldberg-Sachs theorem, see also [12].
The outline of this article is as follows. In Section 2 it is introduced the
notation adopted throughout the article to deal with spinors. Section 3
defines the pure subspaces and present a theorem on algebraic identities
satisfied by them. Then, Section 4 extends the formalism from vector spaces to
fiber bundles over differential manifolds. In addition, this section reviews
the issue of introducing a connection on the spinorial bundle and explore the
freedom in its choice. Section 5 presents some differential results regarding
the pure subspaces, connecting the integrability of an isotropic distribution
with differential constraints satisfied by the elements of the associated pure
subspace. Finally, in Section 6 it is proved a theorem relating twistors to
the existence of integrable maximally isotropic distributions. All over this
work it is being assumed that the vector spaces and the manifolds are even-
dimensional and endowed with a non-degenerate metric of arbitrary signature.
Unless otherwise stated, Einstein summation convention is used, meaning that
repeated indices are summed. The results are all local.
## 2 Spinors and the Space $V\oplus V^{*}$
Given a vector space $\mathcal{V}$ endowed with a non-degenerate inner product
$\langle\,,\rangle$, the Clifford Algebra $Cl(\mathcal{V})$ is an algebra in
this vector space such that:
$v\,u\,+\,u\,v\,=\,2\,\langle
v,u\rangle\,\;\;\;\forall\;v,u\,\in\;\mathcal{V}.$ (1)
The space of spinors associated to $(\,\mathcal{V},\langle\,,\rangle\,)$ is a
vectorial space $S$ where an irreducible and faithful representation of
$Cl(\mathcal{V})$ acts. In even dimensions it is always possible to find a
matrix representation for $Cl(\mathcal{V})$, if $\dim(\mathcal{V})=2n$ then
the least-dimensional faithful representation of this algebra is provided by
$2^{n}\times 2^{n}$ matrices. Therefore, in this case spinors are represented
by the column vectors on which these matrices act, so that the space of
spinors has dimension $2^{n}$.
Now, let us deal with the vector space $V\oplus V^{*}$, where $V$ is an
$n$-dimensional vector space and $V^{*}$ is its dual. The space $V\oplus
V^{*}$ is naturally endowed with a non-degenerate symmetric inner product
$\langle\,,\rangle$ defined by:
$\langle
e+\theta,e^{\prime}+\theta^{\prime}\rangle\,\equiv\,\frac{1}{2}\left[\,\theta(e^{\prime})\,+\,\theta^{\prime}(e)\,\right]\;\;\;\;\;\;e,e^{\prime}\in
V\;\textrm{ and }\;\theta,\theta^{\prime}\in V^{*}\,.$ (2)
Let us introduce the basis $\\{e_{1},e_{2},...,e_{n}\\}$ for $V$ and denote by
$\\{\theta^{1},\theta^{2},...,\theta^{n}\\}$ the dual basis,
$\theta^{i}(e_{j})=\delta_{j}^{i}$. So, thanks to (2) we have that the
following relations hold:
$\langle e_{i},e_{j}\rangle\,=\,0\,,\quad\;\quad\langle
e_{i},\theta^{j}\rangle\,=\,\frac{1}{2}\,\delta_{i}^{j}\,,\quad\;\quad\langle\theta^{i},\theta^{j}\rangle\,=\,0\,.$
The interesting feature of the space $V\oplus V^{*}$ endowed with this natural
inner product is that the spinors can be constructed quite easily. Indeed, the
algebra $Cl(V\oplus V^{*})$ admits a representation in the space $\wedge
V^{*}$, the exterior algebra of $V^{*}$. The action of $V\oplus V^{*}$ vectors
in this representation being defined by:
$(e+\theta)\cdot\varphi\,\equiv\,e\lrcorner\varphi\,+\,\theta\wedge\varphi\quad\quad\;e\in
V\;,\;\theta\in V^{*}\;,\;\varphi\in\wedge V^{*}\,,$ (3)
where $e\lrcorner\varphi$ means the interior product of the vector $e$ on the
form $\varphi$. The action of a scalar is defined in the trivial way,
$\lambda\cdot\varphi=\lambda\varphi$. In order to see that this is indeed a
representation of the Clifford algebra, note that by successive application of
some vector $v=e+\theta$ on the spinor $\varphi$ we get
$\displaystyle v\cdot\left[v\cdot\varphi\right]\,$
$\displaystyle=\,(e+\theta)\cdot\left[e\lrcorner\varphi\,+\,\theta\wedge\varphi\right]$
$\displaystyle=\,e\lrcorner
e\lrcorner\varphi\,+\,e\lrcorner(\theta\wedge\varphi)\,+\,\theta\wedge(e\lrcorner\varphi)\,+\,\theta\wedge\theta\wedge\varphi$
$\displaystyle=\,\theta(e)\,\varphi\,=\,\langle
v,v\rangle\,\varphi\,=\,(v\,v)\cdot\varphi\,,$
where in the last step it was used Eq. (1). Since the space $\wedge V^{*}$ has
$2^{n}$ dimensions we conclude that this is, indeed, the space of spinors,
$S=\wedge V^{*}$ [13].
To simplify the notation, it will be used the abbreviation
$\theta^{12...k}\equiv\theta^{1}\wedge\theta^{2}\wedge...\wedge\theta^{k}$.
Furthermore, in order to avoid any confusion it is worth distinguishing the
elements of $Cl(V{\oplus}V^{*})$ from the spinors of $S=\wedge V^{*}$. With
this intent from now on we shall always denote the spinors by a Greek letter
with a hat on top. For instance, $\theta^{i}$ is an element of $V\oplus V^{*}$
while $\hat{\theta}^{i}$ is a spinor.
In the present formalism an inner product can be easily introduced on the
space of spinors. Such product is non-degenerate and defined up to an
arbitrary scale factor. This scale can be fixed by an arbitrary choice of
basis for $S$, here let us choose the basis
$\\{\hat{1},\hat{\theta}^{i},\hat{\theta}^{ij},...,\hat{\theta}^{12...n}\\}$.
Then the inner product $(\,,):S\times S\rightarrow\mathbb{C}$ is defined to be
such that:
$<\,\hat{\varphi}^{\,t}\wedge\hat{\psi}\,>_{n}\,\,=\,\lambda\,\hat{\theta}^{12...n}\;\;\;\Longleftrightarrow\;\;\;(\hat{\varphi},\hat{\psi})\,\equiv\,\lambda\,.$
(4)
Where $<\hat{\varphi}>_{k}$ means the component of degree $k$ of the form
$\hat{\varphi}$, while $\hat{\varphi}^{\,t}$ means the reverse of
$\hat{\varphi}$. For instance, if
$\hat{\varphi}=\hat{1}+\hat{\theta}^{12...k}$ then
$<\hat{\varphi}>_{k}=\hat{\theta}^{12...k}$ and
$\hat{\varphi}^{\,t}=\hat{1}+\hat{\theta}^{k...21}$. It is not difficult
convincing ourselves that
$(v\cdot\hat{\varphi},\hat{\psi})=(\hat{\varphi},v\cdot\hat{\psi})$ for any
$v\in V\oplus V^{*}$. Such property implies that this inner product is
invariant under the group $Spin_{+}(V\oplus V^{*})$, the double cover of
$SO_{+}(V\oplus V^{*})$. More generally, if $\Omega\in Cl(V\oplus V^{*})$ then
it can be proved that
$(\Omega\cdot\hat{\varphi},\hat{\psi})=(\hat{\varphi},\Omega^{t}\cdot\hat{\psi})$.
If $v_{1},\,\ldots,v_{p}$ are vectors in $V\oplus V^{*}$ then we shall define
$v_{1}\wedge\ldots\wedge v_{p}$ to be the complete antisymmetric part of the
Clifford product $v_{1}v_{2}\cdots v_{p}$. In particular, we have that
$v_{1}\wedge v_{2}=\frac{1}{2!}(v_{1}v_{2}-v_{2}v_{1})$. Then, the so-called
pseudo-scalar of the Clifford algebra $Cl(V\oplus V^{*})$ is defined by:
$\mathcal{I}\,=\,2^{n}\,e_{1}\wedge\theta^{1}\wedge
e_{2}\wedge\theta^{2}\wedge\ldots\wedge e_{n}\wedge\theta^{n}\,.$
It is simple matter to prove that $\mathcal{I}\,\mathcal{I}=1$. So that such
object can be used to split the spinor space as $S=S^{+}\oplus S^{-}$, where
$\hat{\varphi}\,\in\,S^{\pm}$ if
$\mathcal{I}\cdot\hat{\varphi}=\pm\hat{\varphi}$. The elements of $S^{+}$ are
called the Weyl spinors of positive chirality, while those of $S^{-}$ are the
Weyl spinors of negative chirality. It can be shown that $S^{+}$ is spanned by
the spinors of even degree, while $S^{-}$ is spanned by the spinors of odd
degree, so that $\dim(S^{\pm})=2^{n-1}$. More about the spinors of $Cl(V\oplus
V^{*})$ can be found in [13], where this formalism is applied to the space
$TM\oplus TM^{*}$, with $TM$ being the tangent bundle of a manifold $M$. In
order for the reader to get acquainted with the language introduced so far,
let us work out a simple example.
Example 1
When $\dim(V)=2$ we have that $\\{e_{1},e_{2},\theta^{1},\theta^{2}\\}$
provides a basis for $V\oplus V^{*}$. Thus, a basis for the Clifford algebra
$Cl(V\oplus V^{*})$ is formed by the following 16 elements:
$\\{1,\,e_{i},\,\theta^{i},\,e_{1}\wedge
e_{2},\,e_{i}\wedge\theta^{j},\,\theta^{1}\wedge\theta^{2},\,e_{i}\wedge\theta^{1}\wedge\theta^{2},\,e_{1}\wedge
e_{2}\wedge\theta^{i},\,e_{1}\wedge
e_{2}\wedge\theta^{1}\wedge\theta^{2}\\}\,.$
While a basis for the spinor space is given by
$\\{\hat{1},\hat{\theta}^{1},\hat{\theta}^{2},\hat{\theta}^{12}\\}$. Since
$\hat{1}$ and $\hat{\theta}^{12}$ are forms of even degree they are Weyl
spinors of positive chirality, while $\hat{\theta}^{1}$ and $\hat{\theta}^{2}$
have negative chirality. Note, for instance, that the following relations
hold:
$e_{i}\cdot\hat{1}=0\,,\quad\;e_{i}\cdot\hat{\theta}^{j}\,=\,\delta^{j}_{i}\,\hat{1}\,,\quad\;e_{1}\cdot\hat{\theta}^{12}\,=\,\hat{\theta}^{2}\,,\quad\;\theta^{i}\cdot\hat{1}=\hat{\theta}^{i}\,,\quad\;\theta^{1}\cdot\hat{\theta}^{2}=\hat{\theta}^{12}\,,\quad\;\theta^{1}\cdot\hat{\theta}^{12}=0\,.$
Moreover, the non-zero inner products are given by:
$(\hat{1},\hat{\theta}^{12})\,=\,-(\hat{\theta}^{12},\hat{1})\,=\,1\quad\textrm{and}\quad(\hat{\theta}^{1},\hat{\theta}^{2})\,=\,-(\hat{\theta}^{2},\hat{\theta}^{1})\,=\,1\,.$
In particular, note that this inner product is skew-symmetric. $\Box$
Although it may appear too restrictive working with vector spaces of the form
$V{\oplus}V^{*}$, this is not the case at all. Every even-dimensional vector
space endowed with a non-degenerate metric can easily be cast in the form
$V\oplus V^{*}$ when complexified. For example, in the Minkowski space
$\mathbb{R}^{1,3}$ a complex null tetrad basis can always be introduced. In
the standard notation of General Relativity this null tetrad is denoted by
$\\{l,m,n,\overline{m}\\}$, with $\langle l,n\rangle=1$ and $\langle
m,\overline{m}\rangle=-1$, all other inner products between the basis vectors
being zero. So it is possible to make the following associations:
$l\leftrightarrow e_{1}$, $m\leftrightarrow e_{2}$, $n\leftrightarrow
2\theta^{1}$, $\overline{m}\leftrightarrow-2\theta^{2}$. Thus, $\\{l,m\\}$ can
be seen as a basis for $V$ and $\\{n,\overline{m}\\}$ a basis for $V^{*}$. So,
from now on all calculations will be done on vector spaces of the form
$V\oplus V^{*}$ and over the complex field. But it must be clear that the
results can be easily carried to all complexified even-dimensional spaces.
When the metric on the even-dimensional space has split signature the results
are also valid without complexification, since in this case the isotropic
subspaces can have dimension equal to half of the dimension of the full vector
space. The results on real even-dimensional vector spaces can be extracted
from the complex case by choosing suitable reality conditions, in the spirit
of [9, 12].
## 3 Pure Spinorial Subspaces, Algebraic Results
While dealing with spinors it is common, and often valuable, associating to
every non-zero spinor $\hat{\varphi}$ a vector subspace
$N_{\hat{\varphi}}\subset(V\oplus V^{*})$ defined by:
$N_{\hat{\varphi}}=\\{\,v\in(V\oplus V^{*})\;|\;v\cdot\hat{\varphi}=0\,\\}\,.$
These vector subspaces are called isotropic or totally null, because every
vector belonging to $N_{\hat{\varphi}}$ has zero norm. This can be easily
verified, if $v\in N_{\hat{\varphi}}$ then $\langle
v,v\rangle\hat{\varphi}=(vv)\cdot\hat{\varphi}=v\cdot(v\cdot\hat{\varphi})=0$,
so that $\langle v,v\rangle=0$. The maximum dimension that an isotropic
subspace of $V\oplus V^{*}$ can have is $n=\dim(V)$. Therefore, an isotropic
subspace of dimension $n$ is called a maximally isotropic subspace. In the
particular case of $N_{\hat{\varphi}}$ being maximally isotropic, the spinor
$\hat{\varphi}$ is said to be a pure spinor. For example, in the notation of
last section the spinor $\hat{1}$ is such that
$N_{\hat{1}}=\textrm{Span}\\{e_{1},e_{2},...,e_{n}\\}$, so that $\hat{1}$ is a
pure spinor. The pure spinors are very special objects in mathematics and have
been studied since the beginning of XX century, more about them can be found
in [14, 15, 16]. Pure spinors are also acquiring increasing relevance in
physics, particularly in string theory [5].
However, it seems to have been overlooked that it is also possible to go the
other way around and associate spinors to isotropic subspaces. More precisely,
given an isotropic subspace $I$, one can define a spinorial subspace
$\hat{L}_{I}$ spanned by all spinors annihilated by the Clifford action of
$I$. The aim of the present article is to explore this new path.
Suppose that $I\subset(V\oplus V^{*})$ is an isotropic subspace. Then, let us
define the subspace $\hat{L}_{I}\subset S$ as follows:
$\hat{L}_{I}\,\equiv\,\left\\{\,\hat{\varphi}\in
S\;|\;v\cdot\hat{\varphi}=0\;\;\forall\;v\in I\,\right\\}\,.$ (5)
It is trivial to see that $\hat{L}_{I}$ is, indeed, a vector subspace of the
spinor space $S=\wedge V^{*}$. In the particular case of $I$ being maximal,
$\dim(I)=n$, it follows that $\hat{L}_{I}=\textrm{Span}\\{\hat{\psi}\\}$,
where $\hat{\psi}$ is the pure spinor associated to $I$ and
$\textrm{Span}\\{\hat{\psi}\\}$ is the one-dimensional subspace spanned by it.
This section is devoted to enunciate and prove several algebraic results
concerning the spaces $\hat{L}_{I}$, where $I$ is any totaly null subspace of
$V\oplus V^{*}$. Before proceeding let us see a simple example.
Example 2
If $\dim(V)=3$ we have that
$\\{e_{1},e_{2},e_{3},\theta^{1},\theta^{2},\theta^{3},\\}$ is a basis to
$V\oplus V^{*}$. The spinor space is generated by the basis
$\\{\hat{1},\hat{\theta}^{1},\hat{\theta}^{2},\hat{\theta}^{3},\hat{\theta}^{12},\hat{\theta}^{13},\hat{\theta}^{23},\hat{\theta}^{123}\\}$.
Then, let us define the following isotropic subspaces of $V\oplus V^{*}$:
$I_{1}\,=\,\textrm{Span}\\{e_{1}\\}\,,\quad
I_{2}\,=\,\textrm{Span}\\{e_{1},\theta^{2}\\}\,,\quad
I_{3}\,=\,\textrm{Span}\\{e_{1},\theta^{2},\theta^{3}\\}\,.$
Thus, the associated spinorial subspaces are respectively given by:
$\hat{L}_{I_{1}}\,=\,\textrm{Span}\\{\hat{1},\hat{\theta}^{2},\hat{\theta}^{3},\hat{\theta}^{23}\\}\,,\quad\hat{L}_{I_{2}}\,=\,\textrm{Span}\\{\hat{\theta}^{2},\hat{\theta}^{23}\\}\,,\quad\hat{L}_{I_{3}}\,=\,\textrm{Span}\\{\hat{\theta}^{23}\\}\,.$
$\Box$
Now, let us enunciate, in the form of a theorem, several interesting algebraic
results regarding the spinorial subspaces $\hat{L}_{I}$.
###### Theorem 1
Given two isotropic subspaces $I,I^{\prime}\subset(V\oplus V^{*})$, with $V$
being a complexified vector space of dimension $n$, then it follows that:
(1) $\dim(\hat{L}_{I})\,=\,2^{n-\dim(I)}$
(2) $\hat{L}_{I}=\hat{L}_{I^{\prime}}$ $\Longleftrightarrow$ $I=I^{\prime}$
(3) $(\hat{L}_{I}+\hat{L}_{I^{\prime}})\subset\hat{L}_{I\cap I^{\prime}}$
(4) If $I^{\prime\prime}=I+I^{\prime}$ is isotropic then
$\hat{L}_{I+I^{\prime}}=\hat{L}_{I}\cap\hat{L}_{I^{\prime}}$
(5) If $u\in(V\oplus V^{*})$ is such that $u\cdot\hat{\varphi}=0$
$\,\forall\,\;\hat{\varphi}\in\hat{L}_{I}$ then $u\in I$
(6) $I^{\prime}\subset I$ $\Longleftrightarrow$
$\hat{L}_{I}\subset\hat{L}_{I^{\prime}}$
(7) If $I\neq\\{0\\}$ then $(\hat{\varphi},\hat{\psi})=0$
$\,\forall\,\;\hat{\varphi},\hat{\psi}\in\hat{L}_{I}$
(8) $I^{\prime\prime}=I+I^{\prime}$ is isotropic $\Longleftrightarrow$
$\hat{L}_{I}\cap\hat{L}_{I^{\prime}}\neq\\{0\\}$
(9) If $I$ is not maximal then
$\hat{L}_{I}=\hat{L}^{+}_{I}\oplus\hat{L}^{-}_{I}$ where $\hat{L}^{\pm}_{I}$
are spanned by Weyl spinors of $\pm$ chirality and
$\dim(\hat{L}^{+}_{I})=\dim(\hat{L}^{-}_{I})$.
Before proving the above results it is important to note that given an
isotropic subspace $I\subset V\oplus V^{*}$ of dimension $k$ then it is always
possible to find a vector subspace
$V^{\prime}=\textrm{Span}\\{e_{1},\ldots,e_{n}\\}$ such that $I$ is spanned by
$\\{e_{1},e_{2},\ldots,e_{k}\\}$ and $V\oplus V^{*}=V^{\prime}\oplus
V^{\prime*}$. Hence, the judicious and convenient choice
$I=\textrm{Span}\\{e_{1},\ldots,e_{k}\\}$ represents no loss of generality.
_Proof of Theorem 1:_
_(1)_ If $I=\textrm{Span}\\{e_{1},e_{2},\ldots,e_{k}\\}$ then it is not hard
to conclude that
$\hat{L}_{I}=\textrm{Span}\\{\,\hat{1},\,\hat{\theta}^{\alpha^{\prime}},\,\hat{\theta}^{\alpha^{\prime}\beta^{\prime}},\,\cdots,\,\hat{\theta}^{k+1\,k+2...n}\,\\}\,,$
where $\alpha^{\prime},\beta^{\prime}\,\in\\{k+1,k+2,...,n\\}$. Leading us to
the following result:
$\dim(\hat{L}_{I})=1+(n-k)+\frac{(n-k)(n-k-1)}{2!}+\cdots+1=\sum_{i=0}^{n-k}\,\dbinom{n-k}{i}=(1+1)^{n-k}.$
_(2)_ Suppose that $\hat{L}_{I}=\hat{L}_{I^{\prime}}$ but $I\neq I^{\prime}$,
then there exists $u\in I$ such that $u$ does not belong to $I^{\prime}$ (or
the converse). Now, if $\hat{\varphi}\in\hat{L}_{I}=\hat{L}_{I^{\prime}}$ is a
non-zero spinor then $2\langle
u,v^{\prime}\rangle\hat{\varphi}=(uv^{\prime}+v^{\prime}u)\cdot\hat{\varphi}=0$
$\,\forall\,\,v^{\prime}\in I^{\prime}$, so that $\langle
u,v^{\prime}\rangle=0$ for all $v^{\prime}\in I^{\prime}$. This implies that
$I^{\prime\prime}\equiv I^{\prime}+\textrm{Span}\\{u\\}$ is an isotropic
subspace. Moreover, since $u\notin I^{\prime}$ it follows that
$\dim(I^{\prime\prime})=\dim(I^{\prime})+1$, from which we conclude that
$\dim(\hat{L}_{I^{\prime\prime}})<\dim(\hat{L}_{I^{\prime}})$. On the other
hand, since $\hat{L}_{I}=\hat{L}_{I^{\prime}}$ then
$u\cdot\hat{L}_{I^{\prime}}=0$, so if $\hat{\varphi}\in\hat{L}_{I^{\prime}}$
then $v^{\prime\prime}\cdot\hat{\varphi}=0$ for all $v^{\prime\prime}\in
I^{\prime\prime}$. Thus,
$\hat{L}_{I^{\prime}}\subset\hat{L}_{I^{\prime\prime}}$ which implies
$\dim(\hat{L}_{I^{\prime\prime}})\geq\dim(\hat{L}_{I^{\prime}})$,
contradicting the former inequality. Therefore, if
$\hat{L}_{I}=\hat{L}_{I^{\prime}}$ then $I=I^{\prime}$. The converse is
trivial.
_(3)_ If $\hat{\psi}\in(\hat{L}_{I}+\hat{L}_{I^{\prime}})$ then
$\hat{\psi}=\hat{\varphi}+\hat{\varphi}^{\prime}$, where
$v\cdot\hat{\varphi}=0=v^{\prime}\cdot\hat{\varphi}^{\prime}$ for all $v\in I$
and $v^{\prime}\in I^{\prime}$. Thus, if $u\in I\cap I^{\prime}$ then
$u\cdot\hat{\varphi}=0=u\cdot\hat{\varphi}^{\prime}$, which implies
$u\cdot\hat{\psi}=0$. This means that $\hat{\psi}\in\hat{L}_{I\cap
I^{\prime}}$, proving the wanted relation.
_(4)_ Suppose that $I^{\prime\prime}=I+I^{\prime}$ is an isotropic subspace
and let $v^{\prime\prime}\in I^{\prime\prime}$, then
$v^{\prime\prime}=v+v^{\prime}$ where
$v\cdot\hat{\varphi}=0=v^{\prime}\cdot\hat{\varphi}^{\prime}$ for all
$\hat{\varphi}\in\hat{L}_{I}$ and
$\hat{\varphi}^{\prime}\in\hat{L}_{I^{\prime}}$. Thus, if
$\hat{\psi}\in\hat{L}_{I}\cap\hat{L}_{I^{\prime}}$ then
$v^{\prime\prime}\cdot\hat{\psi}=0$. This means that
$\hat{L}_{I}\cap\hat{L}_{I^{\prime}}\subset\hat{L}_{I^{\prime\prime}}$.
Conversely, if $\hat{\psi}\in\hat{L}_{I^{\prime\prime}}$ then
$(v+v^{\prime})\cdot\hat{\psi}=0$ for all $v\in I$ and $v^{\prime}\in
I^{\prime}$. In particular, if we take $v^{\prime}=0$ then we get
$v\cdot\hat{\psi}=0$ for all $v\in I$ $\Rightarrow$
$\hat{\psi}\in\hat{L}_{I}$. In an analogous fashion it is obtained that
$\hat{\psi}\in\hat{L}_{I^{\prime}}$, so that
$\hat{\psi}\in\hat{L}_{I}\cap\hat{L}_{I^{\prime}}$. Therefore, we get
$\hat{L}_{I^{\prime\prime}}\subset\hat{L}_{I}\cap\hat{L}_{I^{\prime}}$,
establishing the equality
$\hat{L}_{I^{\prime\prime}}=\hat{L}_{I}\cap\hat{L}_{I^{\prime}}$.
_(5)_ Suppose that $u\in(V\oplus V^{*})$ is such that $u\cdot\hat{\varphi}=0$
for all $\hat{\varphi}\in\hat{L}_{I}$, then if $v\in I$ it follows that
$2\langle u,v\rangle\hat{\varphi}=(uv+vu)\cdot\hat{\varphi}=0$. Hence,
$\langle u,v\rangle=0$ for all $v$, so that
$I^{\prime\prime}=I+\textrm{Span}\\{u\\}$ is an isotropic subspace of $V\oplus
V^{*}$. By hypothesis it follows that
$\hat{L}_{I}\subset\hat{L}_{I^{\prime\prime}}$. On the other hand, by
definition, we have $\dim(I^{\prime\prime})\geq\dim(I)$ which implies, by part
_(1)_ of this theorem, that
$\dim(\hat{L}_{I^{\prime\prime}})\leq\dim(\hat{L}_{I})$. Therefore, we
conclude that $\hat{L}_{I}=\hat{L}_{I^{\prime\prime}}$. Now, using part _(2)_
of this theorem we get $I=I^{\prime\prime}\equiv I+\textrm{Span}\\{u\\}$, thus
$u\in I$.
_(6)_ If $I^{\prime}\subset I$ then $v^{\prime}\cdot\hat{\varphi}=0$ for all
$v^{\prime}\in I^{\prime}$ and $\hat{\varphi}\in\hat{L}_{I}$, which is
tantamount to $\hat{L}_{I}\subset\hat{L}_{I^{\prime}}$. Conversely, if
$\hat{L}_{I}\subset\hat{L}_{I^{\prime}}$ then $v^{\prime}\cdot\hat{\varphi}=0$
for all $\hat{\varphi}\in\hat{L}_{I}$ and $v^{\prime}\in I^{\prime}$, which
thanks to the item _(5)_ of this theorem, implies that $v^{\prime}\in I$, and
so $I^{\prime}\subset I$.
_(7)_ As explained above, if $\dim(I)=k>0$ then one can assume, without loss
of generality, that $I=\textrm{Span}\\{e_{1},e_{2},\ldots,e_{k}\\}$. With this
assumption $\hat{L}_{I}$ is generated by the basis
$\\{\hat{1},\hat{\theta}^{\alpha^{\prime}},\hat{\theta}^{\alpha^{\prime}\beta^{\prime}},\ldots,\hat{\theta}^{k+1\,k+2...n}\\}$,
where $\alpha^{\prime},\beta^{\prime}\,\in\\{k+1,k+2,...,n\\}$. The inner
product of two elements of this basis,
$(\hat{\varphi}_{1},\hat{\varphi}_{2})$, is always zero, since
$<\hat{\varphi}_{1}^{\,t}\wedge\hat{\varphi}_{2}>_{n}=\pm<\hat{\varphi}_{1}\wedge\hat{\varphi}_{2}>_{n}=0$.
With the last equality stemming from the fact that neither $\hat{\varphi}_{1}$
nor $\hat{\varphi}_{2}$ contains the term $\hat{\theta}^{1}$, so that it is
impossible to get the spinor $\hat{\theta}^{12...n}$ after the wedge product.
Then, using (4) we see that $(\hat{\varphi}_{1},\hat{\varphi}_{2})=0$. The
rest of the statement follows from the bilinearity of this inner product.
_(8)_ If $I^{\prime\prime}=I+I^{\prime}$ is isotropic then, from parts _(1)_
and _(4)_ of this theorem, it follows that
$\dim(\hat{L}_{I}\cap\hat{L}_{I^{\prime}})=\dim(\hat{L}_{I^{\prime\prime}})\geq
1$, so $\hat{L}_{I}\cap\hat{L}_{I^{\prime}}\neq\\{0\\}$. Conversely, if
$\hat{L}_{I}\cap\hat{L}_{I^{\prime}}\neq\\{0\\}$ then there exists
$\hat{\varphi}\neq 0$ such that
$v\cdot\hat{\varphi}=0=v^{\prime}\cdot\hat{\varphi}$ for all $v\in I$ and
$v^{\prime}\in I^{\prime}$. Thus, $2\langle
v,v^{\prime}\rangle\hat{\varphi}=(vv^{\prime}+v^{\prime}v)\cdot\hat{\varphi}=0$,
so that $\langle v,v^{\prime}\rangle=0$ for all $v\in I$ and $v^{\prime}\in
I^{\prime}$, implying that $I^{\prime\prime}=I+I^{\prime}$ is isotropic.
_(9)_ This result is a simple consequence of the fact that if $I$ is non-
maximal, $k\neq n$, then it can always be chosen to be
$I=\textrm{Span}\\{e_{1},e_{2},\ldots,e_{k}\\}$, so that
$\hat{L}_{I}=\textrm{Span}\\{\hat{1},\hat{\theta}^{\alpha^{\prime}},\hat{\theta}^{\alpha^{\prime}\beta^{\prime}},\ldots,\hat{\theta}^{k+1\,k+2...n}\\}$,
where $\alpha^{\prime},\beta^{\prime}\,\in\\{k+1,k+2,...,n\\}$. In this basis
for $\hat{L}_{I}$ there are $2^{n-k-1}$ spinors of even grade and $2^{n-k-1}$
spinors of odd grade. $\Box$
A corollary that easily follows from items _(1)_ , _(4)_ and _(8)_ is that if
two isotropic subspaces $I$ and $I^{\prime}$ are such that the dimension of
$\hat{L}_{I}\cap\hat{L}_{I^{\prime}}$ is not even then
$\dim(\hat{L}_{I}\cap\hat{L}_{I^{\prime}})=1$ and $(I+I^{\prime})$ is a
maximally isotropic subspace. Furthermore, part _(3)_ of Theorem 1 states that
$(\hat{L}_{I}+\hat{L}_{I^{\prime}})$ is contained in the space $\hat{L}_{I\cap
I^{\prime}}$. Then, a natural question to be posed is whether this result can
be made stronger. For instance, is it true that both spaces are equal? The
answer is no, generally there are elements of $\hat{L}_{I\cap I^{\prime}}$
that are not contained in $(\hat{L}_{I}+\hat{L}_{I^{\prime}})$. More
precisely, it can be proved that if
$(\hat{L}_{I}+\hat{L}_{I^{\prime}})=\hat{L}_{I\cap I^{\prime}}$ then one of
the special cases must happen: (i) $I\subset I^{\prime}$, (ii)
$I^{\prime}\subset I$ or (iii) $\dim(I)=\dim(I^{\prime})=\dim(I\cap
I^{\prime})+1$ with $(I+I^{\prime})$ non-isotropic.
It is not hard to see that the pseudo-scalar $\mathcal{I}\in Cl(V\oplus
V^{*})$ anti-commutes with every vector $u\in V\oplus V^{*}$,
$\mathcal{I}u=-u\mathcal{I}$. Thus, if $I\subset V\oplus V^{*}$ is some
isotropic subspace and $\hat{\varphi}\in\hat{L}_{I}$ then
$v\cdot(\mathcal{I}\cdot\hat{\varphi})=-\mathcal{I}\cdot(v\cdot\hat{\varphi})=0$
for all $v\in I$. Therefore, the spinor $(\mathcal{I}\cdot\hat{\varphi})$ must
belong to $\hat{L}_{I}$. In particular, if $\hat{\varphi}\in\hat{L}_{I}$ then
the Weyl spinors of positive and negative chirality
$\hat{\varphi}^{\pm}=\frac{1}{2}(\hat{\varphi}\pm\mathcal{I}\cdot\hat{\varphi})$
belong to $\hat{L}_{I}$ as well. This, in turn, implies that $\hat{L}_{I}$ can
be split as the direct sum of a subspace of positive chirality and a subspace
of negative chirality, which agrees with item _(9)_ of Theorem 1.
Particularly, such reasoning entails the well-known result that every pure
spinor must be a Weyl spinor. It is also worth noting that part _(7)_ of this
theorem guarantees that every pure spinor $\hat{\psi}$ must be orthogonal to
itself, $(\hat{\psi},\hat{\psi})=0$.
Another interesting problem that can be posed is the following: Given a proper
subspace of the spinor space, $\Upsilon\subset S$, then how can one know
whether $\Upsilon$ it is related to some isotropic subspace $I\subset V\oplus
V^{*}$ ? More precisely, given $\Upsilon$ how to know if there exists some
isotropic subspace $I$ such that $\Upsilon=\hat{L}_{I}$? The above theorem
provides several clues for the answer: If one of the conditions (i)
$\dim(\Upsilon)\neq 2^{s}$ for some $s\in\\{0,1,...,n-1\\}$, (ii)
$(\hat{\varphi}_{1},\hat{\varphi}_{2})\neq 0$ for some
$\hat{\varphi}_{1},\hat{\varphi}_{2}\,\in\Upsilon$, (iii) $\dim(\Upsilon)=1$
and $\hat{\varphi}\in\Upsilon$ is not a Weyl spinor or (iv)
$\dim(\Upsilon)\neq 1$ and $\Upsilon$ does not admit a basis such that half of
its elements are Weyl spinors of positive chirality while the other half are
Weyl spinors of negative chirality happen then the spinorial subspace
$\Upsilon$ is not related to any isotropic subspace. But the converse is not
true. For instance, if $\Upsilon$ is the space spanned by some
$\hat{\varphi}^{+}\in S^{+}$ then in general $\Upsilon$ is not related to an
isotropic subspace, since it is well-known that although every pure spinor is
a Weyl spinor not all Weyl spinors are pure111If $n\leq 3$ then all Weyl
spinors are pure. But if $n>3$ there are more Weyl spinors than pure spinors..
Then, it is convenient to make the following definition:
Definition: _A proper spinorial subspace $\Upsilon\subset S$ is called pure if
there exists some isotropic subspace $I\subset V\oplus V^{*}$, such that
$\Upsilon=\hat{L}_{I}$. In the latter case we shall say that $\Upsilon$ is the
pure subspace associated to $I$._
Since the constraint that a spinor must obey in order for it to be pure is
quadratic in the spinor [14, 16], then, probably, the algebraic conditions
that a proper spinorial subspace $\Upsilon\subset S$ might obey in order for
it to be a pure subspace might consist of quadratic equations in the spinors
of $\Upsilon$.
As a last comment in this section, note that the space $\hat{L}_{I}$ can be
defined in a different way from the one presented in Eq. (5). When $I$ is
maximal then $\hat{L}_{I}$ is the pure spinor line related to $I$. When $I$ is
not maximal there are several ways to complete $I$ in order to create
maximally isotropic subspaces. Now, let $\\{I_{a}\\}$ be the set of all
isotropic subspaces of dimension $n$ such that $I\subset I_{a}$. Associated to
each maximally isotropic subspace $I_{a}$ there is a pure spinor line spanned
by $\hat{\varphi}_{a}$. Then, the space $\hat{L}_{I}$ is just the one spanned
by the pure spinors $\\{\hat{\varphi}_{a}\\}$.
Remark:_It is not difficult to note that if
$\hat{\varphi}_{1},\hat{\varphi}_{2}\in\hat{L}_{I}$ and $(e+\theta)\in I$ then
$(e+2\theta)\cdot(\hat{\varphi}_{1}\wedge\hat{\varphi}_{2})=0$. In particular,
this implies that if $I^{\prime}=\textrm{Span}\\{e_{1},e_{2},\ldots,e_{k}\\}$
then
$(\hat{\varphi}^{\prime}_{1}\wedge\hat{\varphi}^{\prime}_{2})\in\hat{L}_{I^{\prime}}$
for all
$\hat{\varphi}^{\prime}_{1},\hat{\varphi}^{\prime}_{2}\in\hat{L}_{I^{\prime}}$.
But this is not true in the case of $I^{\prime}$ being an arbitrary isotropic
subspace. So, probably there exists some bilinear operation $\diamond:S\times
S\rightarrow S$ such that given a general isotropic subspace $I$ then
$(\hat{\varphi}_{1}\diamond\hat{\varphi}_{2})$ belongs to $\hat{L}_{I}$ for
all $\hat{\varphi}_{1},\hat{\varphi}_{2}\in\hat{L}_{I}$ and such that
$\diamond$ reduces to $\wedge$ when $I$ is equal to $I^{\prime}$. _
## 4 Manifolds and Spinorial Connections
From now on we are going to work on a $2n$-dimensional manifold $(M,g)$
endowed with a non-degenerate metric $g$. Furthermore, the tangent bundle $TM$
is assumed to be endowed with a torsion-free and metric-compatible derivative,
the Levi-Civita connection. Since in the present article we are concerned only
with local results it follows that we are allowed to identify the complexified
tangent spaces $\mathbb{C}\otimes T_{p}M$, at any point $p\in M$, with a
vector space of the form $V\oplus V^{*}$, so that the results of the previous
sections can be used. More precisely, if $\\{e_{i},\theta^{j}\\}$ is a frame
in a patch of $M$, with $i,j\in\\{1,2,\ldots,n\\}$, such that
$g(e_{i},e_{j})\,=\,0\,,\quad\;g(e_{i},\theta^{j})\,=\,\frac{1}{2}\,\delta^{j}_{i}\,,\quad\;g(\theta^{i},\theta^{j})\,=\,0\,,$
then we shall say that $V$ is the vector bundle spanned by $\\{e_{i}\\}$ while
$V^{*}$ is the dual vector bundle spanned by $\\{\theta^{i}\\}$, with
$\theta^{j}(e_{i})\equiv\delta^{j}_{i}$. So that an equation analogous to (2)
holds:
$g(e+\theta,e^{\prime}+\theta^{\prime})\,\equiv\,\frac{1}{2}\left[\,\theta(e^{\prime})\,+\,\theta^{\prime}(e)\,\right]\;\;\;\;\;\;e,e^{\prime}\in\Gamma(V)\;\textrm{
and }\;\theta,\theta^{\prime}\in\Gamma(V^{*})\,.$
Where $\Gamma(V)$ is the space of sections of the bundle $V$. So, in this
approach a fiber of the spinorial bundle $S$ will be identified with $\wedge
V^{*}$. It is worth remarking that the identifications $\mathbb{C}\otimes
T_{p}M\sim V|_{p}\oplus V^{*}|_{p}$ and $S_{p}\sim\wedge V^{*}|_{p}$ are not
canonical at all, they depend on an arbitrary choice of frame. Generally, the
frame $\\{e_{i},\theta^{j}\\}$ cannot be globally defined (non-parallelizable
manifolds), so that these objects are well-defined just in a local
neighborhood of the arbitrary point $p\in M$. Therefore, the space $V$ might,
actually, be understood as a local trivialization of some bundle, which is
sufficient for our purposes.
Before proceeding, let us establish some notation. It is useful to denote the
frame $\\{e_{i},\theta^{j}\\}$ by $\\{e_{a}\\}$, where the indices
$a,b,\ldots$ run from $1$ to $2n$ and $e_{i+n}\equiv\theta^{i}$. The
components of the metric in this frame are denoted by the symmetric matrix
$g_{ab}\equiv g(e_{a},e_{b})$, while $g^{ab}$ stands for its inverse,
$g^{ab}g_{bc}=\delta^{a}_{c}$. Finally, let $\omega_{ab}^{\phantom{ab}c}$ be
the connection coefficients in this frame:
$\nabla_{a}\,e_{b}\,\equiv\,\nabla_{e_{a}}\,e_{b}\,=\,\omega_{ab}^{\phantom{ab}c}\,e_{c}\,.$
It is customary to raise or lower the indices by means of $g^{ab}$ and
$g_{ab}$ respectively. For instance,
$\omega_{abd}=\omega_{ab}^{\phantom{ab}c}g_{cd}$ and
$\omega_{a}^{\phantom{a}dc}=\omega_{ab}^{\phantom{ab}c}g^{bd}$. Since the
components of the matrix $g_{ab}$ are constant and the connection is metric
compatible, it follows that $\omega_{abc}=-\omega_{acb}$.
Now, let us introduce a connection $\hat{\nabla}_{a}$ on the spinorial bundle
$S$. In order for this connection to be an extension of the Levi-Civita
connection one might impose for it to satisfy the Leibniz rule with respect to
the Clifford action:
$\hat{\nabla}_{a}\,(v\cdot\hat{\varphi})\,=\,(\nabla_{a}\,v)\cdot\hat{\varphi}\,+\,v\cdot\,\hat{\nabla}_{a}\,\hat{\varphi}\quad\;\;\forall\;\;v\in\Gamma(TM)\;,\,\,\hat{\varphi}\in\Gamma(S)\,.$
(6)
Then, let $\\{\hat{\psi}_{\alpha}\\}$ be a local frame of the spinorial bundle
$S$, with $\alpha,\beta,\ldots$ running from $1$ to $2^{n}$, such that the
Clifford action of the frame $\\{e_{a}\\}$ on the spinors
$\hat{\psi}_{\alpha}$ is constant in the relevant patch of $M$. This means
that the following relation holds:
$e_{a}\cdot\hat{\psi}_{\alpha}\,\equiv\,(e_{a})^{\beta}_{\phantom{\beta}\alpha}\,\hat{\psi}_{\beta}\,,\quad\textrm{with
the matrices }\;(e_{a})^{\beta}_{\phantom{\beta}\alpha}\;\textrm{ being
constant}.$ (7)
In physics the matrices $(e_{a})^{\beta}_{\phantom{\beta}\alpha}$ are the so-
called Dirac matrices. For instance, we can assume that the frame
$\\{\hat{\psi}_{\alpha}\\}$ is given by
$\\{\hat{1},\hat{\theta^{i}},\hat{\theta}^{ij},\ldots,\hat{\theta}^{1\ldots
n}\\}$. Then, without loss of generality, define the action of the connection
on this frame of spinors to be:
$\hat{\nabla}_{a}\,\hat{\psi}_{\alpha}\,\equiv\,\Omega_{a}\cdot\hat{\psi}_{\alpha}\,=\,(\Omega_{a})^{\beta}_{\phantom{\beta}\alpha}\,\hat{\psi}_{\beta}\,.$
(8)
Now, we might look for the general expression of $\Omega_{a}$ that satisfies
Eq. (6). Computing $\hat{\nabla}_{a}(e_{b}\cdot\hat{\psi}_{\alpha})$ in one
hand by means of (7) and on the other hand by means of (6) and then equating
both calculations lead us to the following result:
$(e_{b})^{\beta}_{\phantom{\beta}\alpha}\,(\Omega_{a})^{\sigma}_{\phantom{\sigma}\beta}\,\hat{\psi}_{\sigma}\,\,=\,\,\omega_{ab}^{\phantom{ab}c}\,(e_{c})^{\sigma}_{\phantom{\sigma}\alpha}\,\hat{\psi}_{\sigma}\,+\,(\Omega_{a})^{\beta}_{\phantom{\beta}\alpha}\,(e_{b})^{\sigma}_{\phantom{\sigma}\beta}\,\hat{\psi}_{\sigma}\,.$
In terms of Clifford algebra this equation is tantamount to:
$\omega_{ab}^{\phantom{ab}c}\,e_{c}\,=\,\Omega_{a}\,e_{b}\,-\,e_{b}\,\Omega_{a}\,.$
(9)
Such relation makes clear that $\Omega_{a}$ is defined up to the sum of a term
on the centre of the Clifford algebra. But when the dimension is even, as
assumed throughout this paper, the centre of the Clifford algebra is spanned
by the identity operator. Then, the general solution to Eq. (9) is:
$\Omega_{a}\,=\,-\,\frac{1}{4}\omega_{a}^{\phantom{a}bc}\,(e_{b}\wedge
e_{c})\,+\,A_{a}\,=\,-\,\frac{1}{4}\omega_{a}^{\phantom{a}bc}\,e_{b}\,e_{c}\,+\,A_{a}\,.$
(10)
Where $A_{a}$ are the components of an arbitrary $1$-form
$A\in\Gamma(TM^{*})$. Thus, if $v=v^{b}e_{b}$ and
$\hat{\varphi}=\varphi^{\beta}\hat{\psi}_{\beta}$ are general vector and
spinor fields respectively, then Eqs. (9) and (8) imply that:
$\nabla_{a}\,v\,=\,(\partial_{a}v^{b})e_{b}+v^{b}(\Omega_{a}e_{b}-e_{b}\Omega_{a})\quad;\quad\hat{\nabla}_{a}\,\hat{\varphi}\,=\,(\partial_{a}\varphi^{\beta})\hat{\psi}_{\beta}+\varphi^{\beta}\,\Omega_{a}\cdot\hat{\psi}_{\beta}\,\,.$
(11)
A nice review about connections on the spinorial bundle is available in [18],
see also [19].
Note that since the spinorial space is defined to be the space where an
irreducible and minimal representation of the Clifford algebra acts, it
follows that spinors are defined up to a conformal scale. For instance, if the
frame $\\{\hat{\psi}_{\alpha}\\}$ obeys to Eq. (7) then the frame formed by
the spinors $\hat{\psi}^{\prime}_{\alpha}\equiv
e^{-\lambda}\hat{\psi}_{\alpha}$, for some function $\lambda$, also obeys to
the same equation. Likewise, it is worth recalling that the inner product (4)
was also defined up to a global multiplicative factor, so that one can easily
redefine this inner product, $(\,,)\rightarrow\,\prec\,,\succ$, in such a way
that
$(\hat{\psi}_{\alpha},\hat{\psi}_{\beta})=\prec\hat{\psi}^{\prime}_{\alpha},\hat{\psi}^{\prime}_{\beta}\succ$.
Thus one can see the transformation
$\hat{\varphi}\mapsto\hat{\varphi}^{\prime}=e^{-\lambda}\hat{\varphi}$ as an
intrinsic symmetry of the spinorial formalism. However, in order for this
symmetry to be compatible with the spinorial connection we must assume that
the connection introduced in (10) and (11) change in a way that $A_{a}\mapsto
A^{\prime}_{a}\,=\,A_{a}+\partial_{a}\lambda$. More precisely, we shall define
the connection $\hat{\nabla}^{\prime}_{a}$ to be such that
$\hat{\nabla}^{\prime}_{a}\,\hat{\psi}_{\alpha}\,=\,\Omega^{\prime}_{a}\cdot\hat{\psi}_{\alpha}\,=\,(\Omega_{a}+\partial_{a}\lambda)\cdot\hat{\psi}_{\alpha}\,=\,\hat{\nabla}_{a}\,\hat{\psi}_{\alpha}+(\partial_{a}\lambda)\,\hat{\psi}_{\alpha}\,.$
With this definition one arrives at the following relation:
$\hat{\varphi}^{\prime}\,=\,e^{-\lambda}\,\hat{\varphi}\;\textrm{ and
}A^{\prime}_{a}\,=\,A_{a}+\partial_{a}\lambda\;\;\;\Rightarrow\quad\hat{\nabla}^{\prime}_{a}\,\hat{\varphi}^{\prime}\,=\,e^{-\lambda}\hat{\nabla}_{a}\,\hat{\varphi}\,=\,\left(\hat{\nabla}_{a}\,\hat{\varphi}\right)^{\prime}\,.$
Physically, the above equation says that $\hat{\nabla}_{a}$ is the covariant
derivative associated with the scaling symmetry of the spinors.
The curvature of the spinorial bundle is defined by the following action:
$\hat{\mathfrak{R}}_{ab}\,\hat{\psi}\,=\,\left(\hat{\nabla}_{a}\hat{\nabla}_{b}-\hat{\nabla}_{b}\hat{\nabla}_{a}-\hat{\nabla}_{[e_{a},e_{b}]}\right)\,\hat{\psi}\,.$
Where $[e_{a},e_{b}]$ means the Lie bracket of the vector fields $e_{a}$ and
$e_{b}$. Then, using (8) and (10), one can explicitly prove that this
curvature is given by:
$\hat{\mathfrak{R}}_{ab}\,\hat{\psi}\,=\,-\frac{1}{4}\,R_{ab}^{\phantom{ab}cd}\,(e_{c}e_{d})\cdot\hat{\psi}\,+\,F_{ab}\,\hat{\psi}\;\,\,.$
(12)
Where, in the above equation, $R_{ab}^{\phantom{ab}cd}=g^{cf}g^{dh}R_{abfh}$
is the Riemann tensor of the Riemannian manifold $(M,g)$, while $F_{ab}$ are
the components of the exterior derivative of the $1$-form $A$,
$F_{ab}=(dA)_{ab}$. In particular, note that the spinorial curvature is
invariant under the scaling transformation, $A_{a}\mapsto
A^{\prime}_{a}\,=\,A_{a}+\partial_{a}\lambda$.
Two other operators that are worth mentioning are the Dirac and the twistor
operators, defined respectively by:
$\hat{\boldsymbol{D}}\,=\,g^{ab}e_{a}\cdot\hat{\nabla}_{b}\quad\;;\;\quad\hat{\boldsymbol{T}}_{a}\,=\,\hat{\nabla}_{a}\,-\,\frac{1}{2n}\,e_{a}\cdot\hat{\boldsymbol{D}}\,.$
The twistor operator is characterized by the property
$g^{ab}e_{a}\cdot\hat{\boldsymbol{T}}_{b}=0$. We shall say that a spinor
$\hat{\psi}$ is a twistor if it is annihilated by the action of the twistor
operator, $\hat{\boldsymbol{T}}_{a}\hat{\psi}=0$ $\forall\,a$. As an aside,
note that the square of the Dirac operator is given by:
$\hat{\boldsymbol{D}}^{2}\,\hat{\psi}\,=\,\hat{\Box}\,\hat{\psi}\,+\,\frac{1}{2}F^{ab}\,(e_{a}e_{b})\cdot\,\hat{\psi}\,+\,\frac{1}{4}\,R\,\hat{\psi}\,.$
Where $\hat{\Box}$ is the “spinorial Laplacian”,
$\hat{\Box}=\hat{\nabla}^{a}\hat{\nabla}_{a}-\hat{\nabla}_{\nabla^{a}e_{a}}$,
and $R=R_{ab}^{\phantom{ab}ab}$ is the Ricci scalar.
So far, in order to define a connection on the spinorial bundle we have just
imposed that this connection obeys to the Leibniz rule with respect to the
Clifford action, see (6). As a consequence, we have found that such connection
is unique up to an additive $1$-form $A$. It turns out that such freedom can
be fixed once we require the spinorial connection to obey the Leibniz rule
with respect to the inner product on the spinorial bundle, in addition to Eq.
(6). More precisely, if one requires that
$\nabla_{a}\,(\hat{\varphi},\hat{\psi})\,=\,(\hat{\nabla}_{a}\hat{\varphi},\hat{\psi})\,+\,(\hat{\varphi},\hat{\nabla}_{a}\hat{\psi})$
(13)
then the choice of $A_{a}$ in Eq. (10) is unique. For example, if we choose
the spinorial frame $\\{\hat{\psi}_{\alpha}\\}$ to be
$\\{\hat{1},\hat{\theta^{i}},\hat{\theta}^{ij},\ldots,\hat{\theta}^{1\ldots
n}\\}$ then Eq. (13) holds if, and only if, we set $A_{a}=0$. For instance,
this was the choice of connection made by R. Penrose when he introduced the
spinorial formalism in 4-dimensional general relativity [17], in index
notation Eq. (13) means that the symplectic form $\epsilon_{AB}$ is
covariantly constant. For sake of generality, in the following sections it
will not be assumed that Eq. (13) holds, so that the freedom in the choice of
$A_{a}$ can be exploited.
## 5 Pure Subspaces and Integrability
Once we have introduced a connection on the spinorial bundle, one can look for
theorems on the integrability of isotropic distributions and its relation with
the pure subspaces. For instance, let $I$ be some maximally isotropic
distribution of vector fields over $(M,g)$ and $\hat{\psi}$ the associated
pure spinor, meaning that $\hat{\psi}$ is annihilated by $I$. Then, it is
well-known that the distribution $I$ is integrable if, and only if,
$\hat{\nabla}_{X}\hat{\psi}\propto\hat{\psi}$ for all $X\in I$. Using the
formalism introduced in Section 3, it turns out that this result can be seen
as a particular case of a more general result stated here in the form of the
following theorem.
###### Theorem 2
An isotropic distribution of vector fields $I$ is integrable if, and only if,
$X\cdot\hat{\nabla}_{Y}\hat{\varphi}=Y\cdot\hat{\nabla}_{X}\hat{\varphi}$ for
all $X,Y\in I$ and $\hat{\varphi}\in\hat{L}_{I}$.
_Proof of Theorem 2:_
Let $I=I_{k}=\textrm{Span}\\{e_{1},e_{2},...,e_{k}\\}$, with $k$ fixed, and
$\hat{\varphi}\in\hat{L}_{I_{k}}$. Then, $e_{\beta}\cdot\hat{\varphi}=0$,
where $\beta\in\\{1,2,...,k\\}$. Hence,
$\displaystyle 0=$
$\displaystyle\hat{\nabla}_{\alpha}(e_{\beta}\cdot\hat{\varphi})=(\nabla_{\alpha}e_{\beta})\cdot\hat{\varphi}+e_{\beta}\cdot(\hat{\nabla}_{\alpha}\hat{\varphi})\;\;\Rightarrow$
$\displaystyle[e_{\alpha},e_{\beta}]\cdot\hat{\varphi}=e_{\alpha}\cdot(\hat{\nabla}_{\beta}\hat{\varphi})-e_{\beta}\cdot(\hat{\nabla}_{\alpha}\hat{\varphi})\,.$
(14)
Now, if $I$ is integrable then
$[e_{\alpha},e_{\beta}]=f_{\alpha\beta}^{\gamma}e_{\gamma}$, where
$\alpha,\beta,\gamma\in\\{1,2,...,k\\}$, so that
$[e_{\alpha},e_{\beta}]\cdot\hat{\varphi}=0$. Conversely, by item _(5)_ of
Theorem 1, if $[e_{\alpha},e_{\beta}]\cdot\hat{\varphi}=0$ for all
$\hat{\varphi}\in\hat{L}_{I_{k}}$ then $[e_{\alpha},e_{\beta}]\in I$, so $I$
is integrable. Therefore, (14) implies that $I_{k}$ is integrable
$\Leftrightarrow$
$e_{\alpha}\cdot(\nabla_{\beta}\hat{\varphi})=e_{\beta}\cdot(\nabla_{\alpha}\hat{\varphi})$,
for all $\alpha,\beta\in\\{1,2,...,k\\}$ and for all
$\hat{\varphi}\in\hat{L}_{I_{k}}$, proving Theorem 2. $\Box$
Now, since $I_{k}$ is an isotropic distribution it follows that
$g(e_{\alpha},e_{\beta})=0$, which in terms of Clifford algebra means that
$e_{\alpha}e_{\beta}=-e_{\beta}e_{\alpha}$. Then, supposing that $I_{k}$ is
integrable and using the last theorem we have,
$\displaystyle(e_{\alpha}e_{\beta})\cdot\hat{\nabla}_{\gamma}\hat{\varphi}\,$
$\displaystyle=\,-(e_{\beta}e_{\alpha})\cdot\hat{\nabla}_{\gamma}\hat{\varphi}\,=\,-(e_{\beta}e_{\gamma})\cdot\hat{\nabla}_{\alpha}\hat{\varphi}\,=\,(e_{\gamma}e_{\beta})\cdot\hat{\nabla}_{\alpha}\hat{\varphi}$
$\displaystyle=\,(e_{\gamma}e_{\alpha})\cdot\hat{\nabla}_{\beta}\hat{\varphi}\,=\,-(e_{\alpha}e_{\gamma})\cdot\hat{\nabla}_{\beta}\hat{\varphi}\,=\,-(e_{\alpha}e_{\beta})\cdot\hat{\nabla}_{\gamma}\hat{\varphi}\,,$
with $\hat{\varphi}\in\hat{L}_{I_{k}}$ and
$\alpha,\beta,\gamma\in\\{1,2,...,k\\}$. This means that
$v\cdot(e_{\beta}\cdot\hat{\nabla}_{\gamma}\hat{\varphi})=0$ for all $v\in
I_{k}$. Hence, by definition, we conclude that
$(e_{\beta}\cdot\hat{\nabla}_{\gamma}\hat{\varphi})$ belongs to
$\hat{L}_{I_{k}}$, proving the following corollary.
###### Corollary 1
If $I$ is an integrable isotropic distribution then
$[X\cdot(\hat{\nabla}_{Y}\hat{\varphi})]$ belongs to $\hat{L}_{I}$ for all
$X,Y\in I$ and $\hat{\varphi}\in\hat{L}_{I}$.
In the case of $I_{k}$ being maximal we have
$I_{k}=I_{n}=\textrm{Span}\\{e_{1},e_{2},...,e_{n}\\}$, so that
$\hat{L}_{I_{k}}$ is generated by the pure spinor $\hat{1}$. So, the above
corollary implies that if $I_{n}$ is integrable then
$e_{i}\cdot(\hat{\nabla}_{j}\hat{1})\propto\hat{1}$ for all
$i,j\in\\{1,2,...,n\\}$, which lead us to the relation
$\hat{\nabla}_{j}\hat{1}=\lambda_{j}\hat{1}+\kappa_{ji}\hat{\theta}^{i}$. But
the covariant derivative of a chiral spinor cannot change its chirality, so
that we must have $\kappa_{ji}=0$. Hence, if $I$ is a maximally isotropic
distribution that is integrable then the parallel transport of its pure spinor
in a direction tangent to $I$ does not change the direction of the spinor. By
Theorem 2, it is clear that the converse of this result is also valid. This is
a well-known result that in the language of the pure subspaces assumes the
following form.
###### Corollary 2
A maximally isotropic distribution $I$ is integrable if, and only if,
$\hat{\nabla}_{X}\,\hat{\varphi}\propto\hat{\varphi}$ for all $X\in I$ and
$\hat{\varphi}\in\hat{L}_{I}$.
Now, let $I_{1}=\textrm{Span}\\{e_{1}\\}$ be some one-dimensional distribution
generated by the null vector field $e_{1}$, and $\hat{L}_{I_{1}}$ its
associated pure subspace. So, if $\hat{\varphi}\in\hat{L}_{I_{1}}$ then
$(\nabla_{1}e_{1})\cdot\hat{\varphi}=-e_{1}\cdot(\hat{\nabla}_{1}\hat{\varphi})$.
Such relation, along with item _(5)_ of Theorem 1, implies that the vector
field $e_{1}$ is geodesic if, and only if, $(\hat{\nabla}_{1}\hat{\varphi})$
belong to $\hat{L}_{I_{1}}$ for all $\hat{\varphi}\in\hat{L}_{I_{1}}$. This
simple result is just a particular case of a broader theorem concerning the
totally geodesic character of isotropic foliations. Before stating and proving
such theorem, let us recall some properties of totally geodesic submanifolds.
If $M^{\prime}$ is a submanifold of $(M,g)$ then a point $p\in
M^{\prime}\subset M$ is said to be geodesic when every geodesic of $M$ that is
tangent to $M^{\prime}$ at the point $p$ remains in $M^{\prime}$ forever. The
submanifold $M^{\prime}$ is said to be totally geodesic if all its points are
geodesic. As an example, note that a geodesic curve in $M$ is always a one-
dimensional totally geodesic submanifold of $M$. In order to proceed it is
important to explicitly show which restrictions are imposed to the connection
coefficients by the existence of a totally geodesic submanifold. Let
$\\{x^{1},x^{2},...,x^{2n}\\}$ be local coordinates for $M$ in the
neighborhood of $p\in M^{\prime}\subset M$ such that
$\\{x^{1},x^{2},...,x^{k}\\}$ are local coordinates for $M^{\prime}$ in this
neighborhood. Thus, $\\{\partial_{1},\partial_{2},...,\partial_{k}\\}$ spans
the tangent spaces of $M^{\prime}$ near $p$. Now, let $x^{\mu}(t)$ be a
geodesic of $M$ such that $x(0)=p$ and $\frac{dx}{dt}(0)=\partial_{1}|_{p}$.
Then, using the geodesic equation it is trivial to see that near $p$
$\displaystyle x^{\mu}(t)=$ $\displaystyle
x^{\mu}(0)+t\,\delta_{1}^{\phantom{1}\mu}-\frac{1}{2}t^{2}\,\Gamma_{11}^{\phantom{11}\mu}+O(t^{3})\;\;\Rightarrow$
$\displaystyle\frac{dx^{\mu}}{dt}(t)=\delta_{1}^{\phantom{1}\mu}-t\,\Gamma_{11}^{\phantom{11}\mu}+O(t^{2})\,.$
(15)
Where $\Gamma_{\mu\nu}^{\phantom{\mu\nu}\rho}$ is the Christoffel symbol of
the metric $g$ evaluated at $p$. Now, if $p$ is a geodesic point then
$x^{\mu}(t)$ must be a point in $M^{\prime}$ and $\frac{dx^{\mu}}{dt}(t)$ must
be tangent to $M^{\prime}$ for all $t$. In this case Eq. (15) implies that
$\Gamma_{11}^{\phantom{11}\mu}|_{p}=0$ if $\mu>k$. In general, if $p$ is a
geodesic point then $\Gamma_{\alpha\alpha}^{\phantom{11}\mu}|_{p}=0$ (no sum
in $\alpha$) for all $\alpha\in\\{1,2,...,k\\}$ and
$\mu\in\\{k+1,k+2,...,2n\\}$. It is not hard to note that this is also a
sufficient condition. By means of this result along with the Frobenius
theorem, we are led to the following conclusion: _The leaves of an integrable
distribution $I^{\prime}$ are totally geodesic submanifolds if, and only if,
$\nabla_{X}Y$ is tangent to $I^{\prime}$ for all $X,Y$ tangent to
$I^{\prime}$_.
There are some other equivalent ways to characterize an integrable
distribution of totally geodesic leaves. Let $\\{E_{1},E_{2},...,E_{2n}\\}$ be
vector fields that form a frame in $(M,g)$ such that
$\\{E_{1},E_{2},...,E_{k}\\}$ span the leaves of an integrable distribution.
Then, these leaves are totally geodesic submanifolds if, and only if,
$\nabla_{X}(E_{1}\wedge E_{2}\wedge\cdots\wedge E_{k})\propto(E_{1}\wedge
E_{2}\wedge\cdots\wedge
E_{k})\quad\;\forall\;X\,\in\textrm{Span}\\{E_{1},E_{2},...,E_{k}\\}\,.$
Analogously, if $\\{E^{1},E^{2},...,E^{2n}\\}$ is the dual frame of 1-forms,
$E^{a}(E_{b})=\delta_{b}^{\phantom{b}a}$, then the leaves of the integrable
distribution are totally geodesic if, and only if,
$\nabla_{X}(E^{k+1}\wedge\cdots\wedge E^{2n})\propto(E^{k+1}\wedge\cdots\wedge
E^{2n})\quad\;\forall\;X\,\in\textrm{Span}\\{E_{1},E_{2},...,E_{k}\\}$
Finally, it is worth remarking that an embedded Euclidean manifold
$M^{\prime}$ is totally geodesic if, and only if, its second fundamental form
vanishes. Now, we are ready to prove the following theorem:
###### Theorem 3
An isotropic distribution of vector fields $I$ is integrable and its leaves
are totally geodesic submanifolds if, and only if, $(\nabla_{X}\hat{\varphi})$
belongs to $\hat{L}_{I}$ for all $X\in I$ and $\hat{\varphi}\in\hat{L}_{I}$.
_Proof of Theorem 3:_
If $I$ is integrable and generates totally geodesic submanifolds then, by what
was seen above, it follows that $(\nabla_{X}Y)\in I$ for all $X,Y\in I$, so
that
$0=(\nabla_{X}Y)\cdot\hat{\varphi}=-Y\cdot(\hat{\nabla}_{X}\hat{\varphi})\;\;\;\forall\;X,Y\in
I\;\,\textrm{and}\;\,\hat{\varphi}\in\hat{L}_{I}.$
Since $Y$ is any vector field of $I$ then, by definition,
$(\hat{\nabla}_{X}\hat{\varphi})$ belongs to $\hat{L}_{I}$. Conversely, if
$(\hat{\nabla}_{X}\hat{\varphi})\in\hat{L}_{I}$ for all $X\in I$ and
$\hat{\varphi}\in\hat{L}_{I}$ then it follows that
$Y\cdot(\nabla_{X}\hat{\varphi})=0=X\cdot(\hat{\nabla}_{Y}\hat{\varphi})$. So,
by Theorem 2 this implies that the distribution $I$ is integrable. Moreover,
since $0=Y\cdot(\nabla_{X}\hat{\varphi})=-(\nabla_{X}Y)\cdot\hat{\varphi}$ for
all $\hat{\varphi}\in\hat{L}_{I}$ then by item _(5)_ of Theorem 1 we find that
$(\nabla_{X}Y)$ must belong to $I$, implying that the leaves of this isotropic
distribution are totally geodesic submanifolds. $\Box$
Combining this theorem along with Corollary 2 we find that if a maximally
isotropic distribution is integrable then its leaves are totally geodesic, a
known result that was proved in [10]. Particularly, if the signature of the
manifold is Lorentzian and $I$ is a maximally isotropic distribution then
$\dim(I\cap\overline{I})=1$, with $\overline{I}$ denoting the complex
conjugate of the distribution $I$ [20]. Thus, if $I$ is integrable then
$(I\cap\overline{I})$ generates a null geodesic congruence. In four
dimensions, $2n=4$, if the Ricci tensor vanishes then this geodesic congruence
is shear-free and the Weyl tensor is algebraically special [9].
Note that all manipulations of the present section assumed just that the
spinorial connection $\hat{\nabla}_{a}$ obeys to the Leibniz rule with respect
to the Clifford action, meaning that Eq. (6) holds. Thus, in the above results
the $1$-form $A$ of Eq. (10) is arbitrary. In particular, one can use this
freedom and Corollary 2 to prove that: If $\hat{\varphi}$ is a pure spinor
that generates an integrable maximally isotropic distribution then, it is
always possible to choose $A_{a}$ in (10) to be such that
$\hat{\nabla}_{X}\,\hat{\varphi}=0$ for all $X$ tangent to this distribution.
Moreover, by means of Eq. (12), one can verify that this required $1$-form can
be a pure gauge, $A_{a}=\partial_{a}\lambda$ for some function $\lambda$, if,
and only if, $\hat{\mathfrak{R}}_{XY}\,\hat{\varphi}=0$ for all $X,Y$ tangent
to the distribution.
## 6 Twistor Equation and Integrability of Maximally Isotropic Distributions
It is well-known that in four dimensions a pure spinor obeying to the twistor
equation generates an integrable distribution of isotropic planes. The aim of
the present section is to investigate whether an analogous property holds in
higher dimensions. Namely, the following questions are going to be answered:
Does a pure spinor obeying to the twistor equation necessarily generate an
integrable maximally isotropic distribution? What about the converse, does a
pure spinor generating an integrable distribution obeys to the twistor
equation when we judiciously use the freedom in the choice of the spinorial
connection?
If $\hat{\varphi}\neq 0$ is an arbitrary pure spinor then one can always make
a convenient choice of frame such that $\hat{\varphi}=\hat{1}$. In order to
facilitate the calculations, let us assume that this choice was made, meaning
that $e_{i}\cdot\hat{\varphi}=0$ and
$\theta^{i}\cdot\hat{\varphi}=\hat{\theta}^{i}$. Where, as previously defined,
the indices $i,j,k,l$ run from $1$ to $n$ while the indices $a,b$ belong to
$\\{1,2,\ldots,2n\\}$, with $2n$ being the dimension of the manifold. Then,
one can easily prove the following relations:
$\displaystyle(e_{i}\theta^{k}\theta^{l})\cdot\hat{\varphi}=(\delta^{k}_{i}\,\hat{\theta}^{l}-\delta^{l}_{i}\,\hat{\theta}^{k})\quad\;;$
$\displaystyle\;\quad(e_{j}e_{i}\theta^{k}\theta^{l})\cdot\hat{\varphi}=(\delta^{k}_{i}\,\delta^{l}_{j}-\delta^{l}_{i}\,\delta^{k}_{j})\,\hat{1}$
$\displaystyle\hat{\nabla}_{a}\hat{\varphi}=(A_{a}\,+\,\frac{1}{2}\,$
$\displaystyle\omega_{aj}^{\phantom{aj}j})\,\hat{1}+\omega_{aij}\,\hat{\theta}^{ji}\,.$
Where Eqs. (8) and (10) were used. By means of the above relations one can
prove, after some algebra, that the action of the twistor operator in
$\hat{\varphi}$ is given by:
$\displaystyle\hat{\boldsymbol{T}}_{j}\,\hat{\varphi}\,=\,$
$\displaystyle\frac{1}{n}\left[(n-1)(A_{j}+\frac{1}{2}\,\omega_{jk}^{\phantom{jk}k})\,+\,\omega^{i}_{\phantom{i}ij}\right]\hat{1}\,+$
$\displaystyle+\,$
$\displaystyle\frac{1}{n}\left[(n-1)\omega_{jkl}\,+\,2\,\omega_{[lk]j}\right]\,\hat{\theta}^{lk}$
(16) $\displaystyle\hat{\boldsymbol{T}}_{j+n}\,\hat{\varphi}\,=$
$\displaystyle\,\frac{1}{2}\left(A^{j}-\frac{1}{2}\omega^{ji}_{\phantom{ji}i}\right)\,\hat{1}\,+\,\frac{1}{n}\,\omega_{[ikl]}\,\hat{\theta}^{jikl}\,+$
$\displaystyle+\,$
$\displaystyle\left[\,\frac{1}{2}\,\omega^{j}_{\phantom{j}kl}\,+\,\frac{1}{n}\,\delta^{j}_{[k}A_{l]}\,+\,\frac{1}{2n}\,\delta^{j}_{[k}\omega_{l]i}^{\phantom{l]i}i}\,+\,\frac{1}{n}\,\omega^{i}_{\phantom{i}i[k}\delta^{j}_{l]}\,\right]\,\hat{\theta}^{lk}\,.$
(17)
Where in the above equations $V_{[a_{1}a_{2}\ldots a_{p}]}$ means the anti-
symmetric part of $V_{a_{1}a_{2}\ldots a_{p}}$. For instance,
$V_{[ab]}=\frac{1}{2!}(V_{ab}-V_{ba})$. If $\hat{\varphi}$ is a twistor then
the right hand side of both equations above must vanish for all $j$, let us
analyse these two conditions separately. First note that if the dimension is
different from six, $n\neq 3$, then the right hand side of (16) vanishes if,
and only if,
$A_{j}\,=\,\frac{1}{n-1}\,\omega^{i}_{\phantom{i}ji}-\frac{1}{2}\omega_{jk}^{\phantom{jk}k}\quad\;\;\textrm{and}\quad\;\;\omega_{ijk}\,=\,0\,.$
(18)
While in six dimensions $\hat{\boldsymbol{T}}_{j}\hat{\varphi}=0$ if, and only
if,
$A_{j}\,=\,\frac{1}{3-1}\,\omega^{i}_{\phantom{i}ji}-\frac{1}{2}\omega_{jk}^{\phantom{jk}k}\quad\;\;\textrm{and}\quad\;\;\omega_{ijk}\,=\,\omega_{[ijk]}\,.$
Since the maximally isotropic distribution associated to the pure spinor
$\hat{\varphi}=\hat{1}$ is spanned by $\\{e_{1},e_{2},\ldots,e_{n}\\}$ and
this distribution is integrable if, and only if, $\omega_{ijk}=0$ then the
above conditions lead us to the following theorem:
###### Theorem 4
If a pure spinor $\hat{\varphi}$ generates an inetgrable maximally isotropic
distribution then one can always choose the 1-form $A_{a}$ so that
$\hat{\boldsymbol{T}}_{X}\hat{\varphi}=0$ for all vector fields $X$ tangent to
such distribution. Conversely, if the dimension is different from six and
$\hat{\boldsymbol{T}}_{X}\hat{\varphi}=0$ for all $X$ tangent to the maximally
isotropic distribution generated by $\hat{\varphi}$ then this distribution is
integrable.
Now, assume that we choose the $1$-form $A$ so that the components $A_{j}$ are
just as in Eq. (18). Then, inserting this choice into (17) leads us to the
following equation:
$\displaystyle\hat{\boldsymbol{T}}_{j+n}\,\hat{\varphi}\,=$
$\displaystyle\,\frac{1}{2}\left(A^{j}-\frac{1}{2}\omega^{ji}_{\phantom{ji}i}\right)\,\hat{1}\,+\,\frac{1}{n}\,\omega_{[ikl]}\,\hat{\theta}^{jikl}\,+$
$\displaystyle+$
$\displaystyle\,\left[\,\frac{1}{2}\,\omega^{j}_{\phantom{j}kl}\,+\,\frac{1}{n-1}\,\omega^{i}_{\phantom{i}i[k}\delta^{j}_{l]}\,\right]\,\hat{\theta}^{lk}\,.$
(19)
So, in order for the equation $\hat{\boldsymbol{T}}_{j+n}\hat{\varphi}=0$ to
hold the three terms on the right hand side of the above equation must vanish.
The first term can always be made to vanish by a suitable choice of the
$1$-form $A$, namely we must set
$2A_{j+n}\,=\,A^{j}\,=\,\frac{1}{2}\omega^{ji}_{\phantom{ji}i}\,.$
Due to the anti-symmetry in its indices, the spinor $\hat{\theta}^{jikl}$ is
zero in four and six dimensions. Therefore, in these cases the second term on
the right hand side of Eq. (19) is automatically zero. While if the dimension
is greater or equal to eight this term will vanish if, and only if,
$\omega_{[ijk]}=0$. But it is worth recalling that if we assume that
$\hat{\boldsymbol{T}}_{j}\hat{\varphi}=0$ then this condition is already
satisfied. Finally, the third term on the right hand side of Eq. (19) vanish
if, and only if,
$\omega^{i}_{\phantom{i}jk}\,=\,0\;\;\;\forall\;i\neq
j,k\,\quad\textrm{and}\quad\,\omega^{i}_{\phantom{i}ik}\,=\,\omega^{j}_{\phantom{j}jk}\;(\textrm{no
sum in }i,j)\;\;\;\forall\;i\neq k\neq j\,.$ (20)
Where it is worth reemphasizing that in the second condition above no sum is
assumed in the repeated indices $i$ and $j$. Since in four dimensions
$i,j,k\in\\{1,2\\}$ it follows that the conditions of (20) are identically
satisfied in this case. While in higher dimensions these constraints are non-
trivial. In order to give some sort of insight on the meaning of the
constraints displayed in (20), let us register that the integrability of the
maximally isotropic distribution $\textrm{Span}\\{e_{i}\\}$ along with (20) is
tantamount to the following restrictions:
$\displaystyle
g(\nabla_{X}e_{i}-\nabla_{Y}e_{j},e_{k})\,=\,0\quad\;\forall\;\,X,Y\in\,\Gamma(TM)\textrm{
such that}$ (21) $\displaystyle\;g(X,e_{i})\,=\,g(Y,e_{j})\;\textrm{ and
}\;g(X,e_{k})\,=\,g(Y,e_{k})\,=\,0\,.$
The table below summarizes the joint analysis of Eqs. (16) and (17),
displaying the necessary and sufficient conditions for the pure spinor
$\hat{\varphi}$ to be a twistor. Since the condition $\omega_{ijk}=0$ is
equivalent to the integrability of the maximally isotropic distribution
generated by $\hat{\varphi}=\hat{1}$, such table leads us to Theorem 5,
presented on the sequence.
Dimension | $A_{j}$ | $A_{j+n}$ | Constraints over $\omega_{abc}$
---|---|---|---
4 | $\frac{1}{2-1}\omega^{i}_{\phantom{i}ji}-\frac{1}{2}\omega_{jk}^{\phantom{jk}k}$ | $\frac{1}{4}\omega^{ji}_{\phantom{ji}i}$ | $\omega_{ijk}=0$
6 | $\frac{1}{3-1}\omega^{i}_{\phantom{i}ji}-\frac{1}{2}\omega_{jk}^{\phantom{jk}k}$ | $\frac{1}{4}\omega^{ji}_{\phantom{ji}i}$ | $\omega_{ijk}=\omega_{[ijk]}$ , Eq. (20)
$2n\geq 8$ | $\frac{1}{n-1}\omega^{i}_{\phantom{i}ji}-\frac{1}{2}\omega_{jk}^{\phantom{jk}k}$ | $\frac{1}{4}\omega^{ji}_{\phantom{ji}i}$ | $\omega_{ijk}=0$ , Eq. (20)
Table 1: Depending on the dimension, the pure spinor $\hat{\varphi}=\hat{1}$
is a twistor if, and only if, the components of the $1$-form $A$ and the
connection coefficients $\omega_{abc}$ are as displayed in this table. Recall
that the constraints in the last row of the last column can be replaced by the
single Eq. (21).
###### Theorem 5
If $\hat{\varphi}\neq 0$ is a pure spinor then the following results hold
depending on the dimension of the manifold:
(1) $2n=4\rightarrow$ If $\hat{\varphi}$ is a twistor then the maximally
isotropic distribution generated by $\hat{\varphi}$ is integrable. Conversely,
if $\hat{\varphi}$ generates an integrable distribution then one can always
choose $A_{a}$ so that $\hat{\varphi}$ obeys to the twistor equation with
respect to the connection of Eq. (10).
(2) $2n=6\rightarrow$ The fact that $\hat{\varphi}$ obeys to the twistor
equation does not imply that the distribution generated by $\hat{\varphi}$ is
integrable. Conversely, the integrability of the distribution generated by
$\hat{\varphi}$ does not guarantee that one can arrange $A_{a}$ in order to
make $\hat{\varphi}$ obey to the twistor equation.
(3) $2n\geq 8\rightarrow$ If $\hat{\varphi}$ is a twistor then the maximally
isotropic distribution generated by $\hat{\varphi}$ is integrable. On the
other hand, the integrability of the distribution generated by $\hat{\varphi}$
does not imply that one can find $A_{a}$ such that $\hat{\varphi}$ obeys to
the twistor equation.
Hopefully, the results presented in the last two sections will be valuable to
some branches differential geometry. Specially, since the null directions play
an important role on the study of holonomy [21, 22] it follows that some
applications on this subject might appear. Although the study of Euclidean
restricted holonomy is well established [23], for other signatures some
questions remain open. In particular, recently some progress has been
accomplished in the Lorentzian case [22]. It is also worth mentioning the
interplay between holonomy groups and the twistor equation [24]. Concerning
physics, the formalism introduced in the present article might be useful for
string theory and supergravity [25] as well as in the study of geometrical
properties of Black-Holes [26].
## Acknowledgments
I want to thank CNPq (Conselho Nacional de Desenvolvimento Científico e
Tecnológico - Brazil) and CAPES (Coordenação de Aperfeiçoamento de Pessoal de
Nível Superior - Brazil) for the financial support. I also wish to thank the
anonymous referee for helpful comments.
## References
* [1] E. Cartan, Les groupes projectifs qui ne laissent invariante aucune multiplicité plane, Bull. Soc. Math. France 41 (1913), 53.
* [2] W. Graf, Differential forms as spinors, Ann. Inst. Henri Poincaré XXIX (1978), 85; I. Benn and R. Tucker, Fermions without spinors, Commun. Math. Phys. 89 (1983), 341.
* [3] E. Witten, A new proof of the positive energy theorem, Commun. Math. Phys. 80 (1981), 381.
* [4] R. Britto, F. Cachazo, B. Feng and E. Witten, Direct proof of tree-level recursion relation in Yang-Mills theory, Phys. Rev. Lett. 94 (2005), 181602. arXiv:hep-th/0501052
* [5] N. Berkovits and D. Marchioro, Relating the Green-Schwarz and pure spinor formalisms for the superstring, J. High Energy Phys. 01(2005). arXiv:hep-th/0412198
* [6] R. Penrose, A spinor approach to general relativity, Ann. Phys. 10 (1960), 171; E. Newman and R. Penrose, An approach to gravitational radiation by a method of spin coefficients, J. Math. Phys. 3 (1962), 566.
* [7] W. Kinnersley, Type D vacuum metrics, J. Math. Phys. 10 (1969), 1195.
* [8] J. F. Plebański and S. Hacyan, Null geodesic surfaces and Goldberg-Sachs theorem in complex Riemannian spaces, J. Math. Phys. 16 (1975), 2403.
* [9] C. Batista, A generalization of the Goldberg-Sachs theorem and its consequences, Gen. Relativ. Gravit. 45 (2013), 1411. arXiv:1205.4666
* [10] L. Mason and A. Taghavi-Chabert, Killing-Yano tensors and multi-hermitian structures, J. Geom. Phys. 60 (2010), 907. arXiv:0805.3756
* [11] A. Taghavi-Chabert, The complex Goldberg-Sachs theorem in higher dimensions, J. Geom. Phys. 62 (2012), 981. arXiv:1107.2283
* [12] C. Batista, On the Weyl tensor classification in all dimensions and its relation with integrability properties, J. Math. Phys. 54 (2013), 042502. arXiv:1301.2016
* [13] M. Gualtieri, Generalized complex geometry, PhD Thesis, University of Oxford (2003).
* [14] E. Cartan, The theory of spinors, Dover (1966).
* [15] C. Chevalley, The algebraic theory of spinors, Columbia University Press (1954)
* [16] P. Charlton, The geometry of pure spinors, with applications, PhD Thesis, Newcastle University (1997).
* [17] R. Penrose and W. Rindler, Spinors and space-time 1 and 2, Cambridge University Press (1986).
* [18] A. Trautman, Connections and the Dirac operator on spinor bundles, J. Geom. Phys. 58 (2008), 238.
* [19] I. Benn and R. Tucker, An introduction to spinors and geometry with applications in Physics, Adam Hilger (1987).
* [20] W. Kopczynski and A. Trautman, Simple spinors and real structures, J. Math. Phys. 33 (1992), 550.
* [21] H. Baum, K. Lärz and T. Leistner, On the full holonomy group of special Lorentzian manifolds (2012). arXiv:1204.5657
* [22] T. Leistner, On the classification of Lorentzian holonomy groups, J. Differential Geom. 76 (2007), 423. arXiv:math/0305139 ; A. S. Galaev, Metrics that realize all Lorentzian holonomy algebras, Int. J. Geom. Methods Mod. Phys. 3 (2006), 1025. arXiv:math/0502575
* [23] M. Berger, Sur les groupes d’holonomie homogène des variétés à connexion affine et des variétés riemanniennes, Bull. Soc. Math. France 83 (1955), 279.
* [24] A. Lischewski, Towards a classification of pseudo-Riemannian geometries admitting twistor spinors (2013). arXiv:1303.7246
* [25] J. M. Figueroa-O’Farrill, Breaking the M-waves, Class. Quant. Grav. 17 (2000), 2925. arXiv:hep-th/9904124
* [26] B. Carneiro da Cunha and A. R. de Queiroz, Killing Horizons and Spinors (2013). arXiv:1312.4911
|
arxiv-papers
| 2013-10-01T00:19:21 |
2024-09-04T02:49:51.768056
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Carlos Batista",
"submitter": "Carlos A. Batista da S. Filho",
"url": "https://arxiv.org/abs/1310.0104"
}
|
1310.0212
|
[mult]
# MODELING MICROWAVE/ELECTRON-CLOUD INTERACTION
M. Mattes [email protected] E. Sorolla [email protected] EPFL-STI-IEL-
LEMA Station 11 1015 Lausanne Switzerland
F. Zimmermann [email protected] CERN Accelerators and Beam Physics
Group 1211 Geneva 23 Switzerland
###### Abstract
Starting from the separate codes BI-RME and ECLOUD or PyECLOUD, we are
developing a novel joint simulation tool, which models the combined effect of
a charged particle beam and of microwaves on an electron cloud. Possible
applications include the degradation of microwave transmission in tele-
communication satellites by electron clouds; the microwave-transmission
tecchniques being used in particle accelerators for the purpose of electron-
cloud diagnostics; the microwave emission by the electron cloud itself in the
presence of a magnetic field; and the possible suppression of electron-cloud
formation in an accelerator by injecting microwaves of suitable amplitude and
frequency. A few early simulation results are presented.
## Section 1 Motivation
Electron multiplication on surfaces exposed to either an oscillating
electromagnetic field or to a pulsed electric field gives rise to the
phenomenon of multipacting, which can significantly degrade the performance of
radiofrequency devices operating in high-vacuum conditions. These include, for
example, components for satellite tele-communication like waveguide filters
but also accelerating cavities for particle accelerators. In both examples,
the electron multiplication can result in a quasi-stationary “electron cloud”
inside the devices perturbing their performance.
Though studies of electron multipacting in microwave devices as well as of
beam-induced multipactor in particle accelerators have both been ongoing
separately for several decades already [1, 2, 3, 4], little is known about the
simultaneous interaction of an electron cloud with both microwaves and a
particle beam [5, 6]. This lack of knowledge has so far prevented the proper
interpretation of microwave transmission measurements which have been applied
in several accelerators afflicted by an electron cloud. A related phenomenon
is encountered in satellite-based telecommunication systems, which nowadays
suffer from the problem of an electron-cloud build up interacting with the
microwave field used to transmit information.
## Section 2 Model
In the associated oral contribution at ECLOUD12 we have outlined the model
applied to describe the interaction between electrons and microwaves.
Differently from [7], the modelling is based on time-domain Green functions
describing the radiation of a single electron represented as a point charge.
Using the modal expansion representation of the involved Green functions, we
are able to take into account the exact boundary conditions of the devices,
e.g. waveguides or cavities.
## Section 3 Example Results
Figure 1 shows the example of the LHC beam pipe representing the transverse
electric (red arrows) and magnetic field (blue arrows) of the first TE and TM
mode, respectively. The modes are computed using the BI-RME algorithm [8].
Fig. 1: First TE (top) and TM (bottom) mode of LHC beam pipe. The red arrows
represent the transverse electric field, the blue ones the transverse magnetic
field.
In Fig. 2 the example of a single electron under multipactor regime between
two parallel plates separated by 0.2 mm is depicted. The electric field,
driving the electron, oscillates at 5 GHz. The upper figure shows the position
of the electron, the lower one the part of the radiated electric field density
related to the time-derivative of the vector potential ${A}$. The spectrum
includes all harmonics that are odd multiples of the fundamental one and are
generated by the moving electron.
The model developed is useful for microwave diagnosis of devices prone to
microwave breakdown due to multipaction. The upper picture of Fig. 3 shows the
evolution of the electron density for the case of the parallel plate waveguide
of Fig. 2 applying the model of [9]. The blue curve represents the electron
density, the red one the associated radiated electric field contribution from
the time-derivative of the vector potential ${A}$. There is a clear
correlation between the two curves. The lower part of Fig. 3 is an example
from the LHC beam pipe. The space charge build up was simulated with PyECLOUD
[10] (blue curve). Based on the electron trajectories the radiated
electromagnetic field has been computed (red curve). Again, a correlation of
the two curves can be observed.
Fig. 2: Multipacting electron (top) and the associated radiated electric field
spectrum (bottom) inside a parallel-plate waveguide of height 0.2mm.
Simulation parameters: RF frequency: 5 GHz; RF voltage: 70 V, initial velocity
of electrons: 2 eV.
Fig. 3: Example of an electron cloud build-up inside the parallel plate
waveguide of Fig. 2 (top) and in the LHC beam pipe. The red curve repreesents
the electron density (top) respectively the space charge (bottom). The blue
lines are the associated time derivatives of the vector potential${A}$, thus
the electric field.
## Section 4 ACKNOWLEDGMENT
This project has been financed by the Swiss National Science Foundation (SNF)
under Grant No. 200021_129661 “Modelling microwave-electron interaction in the
LHC beam pipe”.
## References
* [1] K. Ohmi, Beam and Photo-electron Interactions in Positron Storage Rings, Phys. Rev. Lett. 75, 1526-1529 (1995)
* [2] M.A. Furman and G.R. Lambertson, The Electron-Cloud Instability in the Arcs of the PEP-II Positron Ring, presented at the International Workshop on Multibunch Instabilities in Future Electron and Positron Accelerators (MBI97), KEK, Tsukuba, Japan, 15-18 July 1997, KEK Proceedings 97-17 (1997) p. 170
* [3] F. Zimmermann, A Simulation Study of Electron-Cloud Instability and Beam-Induced Multipacting in the LHC, CERN LHC Project Report 95 (1997)
* [4] C. Vicente, A. Cervello, M. Mattes, D. Wolk, D. Raboso, B. Mottet, H.L. Hartnagel, and J.R. Mosig, Multipactor and corona discharge prediction in rectangular waveguide based microwave components, in Proc. International Workshop on Microwave Filters, CNES, Toulouse, France, 2004
* [5] F.J. Decker, F. Caspers and F. Zimmermann, Impact of microwaves on the electron cloud and incoherent effects, in Proc. Mini-Workshop on Electron Cloud Simulations for Proton and Positron Beams (ECLOUD’02), pages 277-282, CERN, Geneva, Switzerland, 2002
* [6] F. Caspers and F. Zimmermann, Interactions of microwaves and electron clouds, in 2009 IEEE Particle Accelerator Conference, PAC’09, program code WE1PBI02, Vancouver, B.C., Canada, 2009
* [7] Edén Sorolla et al, An analytical model to evaluate the radiated power spectrum of a multipactor discharge in a parallel-plate region, IEEE Transactions on Electron Devices, 55(8):2252-2258, August 2008
* [8] G. Conciauro, R. Sorrentino and M. Guglielmi, Advanced modal analysis, John Wiley & Sons Inc., New York, NY, USA, 1999d.
* [9] Edén Sorolla and Michael Mattes, Multipactor saturation in parallel-plate waveguides, Physics of Plasmas, Vol. 19, p. 072304, 2012
* [10] G. Iadarola and G. Rumolo, PyECLOUD and build up simulations at CERN, in Proc.of ECLOUD12, La Biodola, Isola d’Elba, Italy, 5-9 June 2012.
|
arxiv-papers
| 2013-10-01T09:43:15 |
2024-09-04T02:49:51.781465
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "M. Mattes (EPFL, Lausanne), E. Sorolla (EPFL, Lausanne), F. Zimmermann\n (CERN)",
"submitter": "Scientific Information Service CERN",
"url": "https://arxiv.org/abs/1310.0212"
}
|
1310.0239
|
# Exact reconstruction formulas for a Radon transform over cones
Markus Haltmeier
###### Abstract
Inversion of Radon transforms is the mathematical foundation of many modern
tomographic imaging modalities. In this paper we study a conical Radon
transform, which is important for computed tomography taking Compton
scattering into account. The conical Radon transform we study integrates a
function in $\mathbb{R}^{d}$ over all conical surfaces having vertices on a
hyperplane and symmetry axis orthogonal to this plane. As the main result we
derive exact reconstruction formulas of the filtered back-projection type for
inverting this transform.
Keywords. Radon transform, conical projections, computed tomography, inversion
formula, filtered back-projection.
AMS subject classifications. 44A12, 45Q05, 92C55.
## 1 Introduction
Suppose that $f\colon\mathbb{R}^{d}\to\mathbb{R}$, with $d\geq 2$, is a smooth
function supported in the half space
$\mathbb{R}^{d-1}\times\left(0,\infty\right)$, and let $p$ be some real
number. We study the problem of reconstructing the function $f$ from the
integrals
$\left(\mathcal{R}^{(p)}f\right)\left(\mathbf{u},\theta\right)\coloneqq\int_{0}^{\infty}\frac{1}{s^{p}}\int_{S^{d-2}}f\left(\mathbf{u}+s\sin\left(\theta\right)\mathbf{n},s\cos\left(\theta\right)\right)\left(s\sin\left(\theta\right)\right)^{d-2}\mathrm{d}\mathbf{n}\mathrm{d}s$
(1.1)
for $\mathbf{u}\in\mathbb{R}^{d-1}$ and $\theta\in\left(0,\pi/2\right)$. (Here
$S^{d-2}$ is the unit sphere in $\mathbb{R}^{d-1}$ and $\mathrm{d}\mathbf{n}$
the surface measure on $S^{d-2}$). We call the function
$\mathcal{R}^{(p)}f\colon\mathbb{R}^{d-1}\times(0,\pi/2)\to\mathbb{R}$ the
conical Radon transform of $f$. As illustrated in Figure 1.1,
$\bigl{(}\mathcal{R}^{(p)}f\bigr{)}\left(\mathbf{u},\theta\right)$ is the
integral of $f$ over the one sided conical surface
$C\left(\mathbf{u},\theta\right)$ having vertex $\left(\mathbf{u},0\right)$ on
the plane $\mathbb{R}^{d-1}\times\left(0,\infty\right)$, symmetry axis
$\mathbf{e}_{d}\coloneqq\left(0,\dots,0,1\right)$, and half opening angle
$\theta\in\left(0,\pi/2\right)$. The product
$\left(s\sin\left(\theta\right)\right)^{d-2}\mathrm{d}\mathbf{n}\mathrm{d}s$
is the standard surface measure on $C\left(\mathbf{u},\theta\right)$, and
$1/s^{p}$ is an additional radial weight that can be adapted to a particular
application at hand. For $\theta\in(0,\pi/2)$, the function
$\bigl{(}\mathcal{R}^{(p)}f\bigr{)}\left(\,\cdot\,,\theta\right)$ may be
considered as a conical projection of $f$ onto
$\mathbb{R}^{d-1}\times\left\\{0\right\\}$.
Figure 1.1: The conical Radon transform integrates a function with support in
the upper half space over one sided conical surfaces
$C\left(\mathbf{u},\theta\right)$ centered at
$(\mathbf{u},0)\in\mathbb{R}^{d-1}\times\left(0,\infty\right)$ having symmetry
axis $\mathbf{e}_{d}=\left(0,\dots,0,1\right)$ and half-angle
$\theta\in\left(0,\pi/2\right)$. Any point on
$C\left(\mathbf{u},\theta\right)$ can be written in the form
$\left(\mathbf{u}+s\sin\left(\theta\right)\mathbf{n},s\cos\left(\theta\right)\right)$
with $\mathbf{n}\in S^{d-2}$ and $s>0$. The $d-1$ dimensional surface measure
on $C\left(\mathbf{u},\theta\right)$ is given by
$\left(s\sin\left(\theta\right)\right)^{d-2}\mathrm{d}\mathbf{n}\mathrm{d}s$,
with $\mathrm{d}\mathbf{n}$ denoting the standard surface measure on
$S^{d-2}$.
Inversion of the conical Radon transform in three spatial dimension is
important for computed tomography taking Compton scattered photons into
account [BZG98, CB94, Par00]. In [CB94, NTG05] Fourier reconstruction formulas
have been derived for the cases $p\in\left\\{0,2\right\\}$. For two spatial
dimensions, $\mathcal{R}^{(p)}$ has been studied with
$p\in\left\\{0,2\right\\}$ in [BZG97, TN11], where reconstruction formulas of
the back-projection type have been derived. In dimensions $d>3$, the conical
Radon transform has, to the best of our knowledge, not been studied so far. In
this paper we study $\mathcal{R}^{(p)}$ for any $d\geq 2$ and any
$p\in\mathbb{R}$. We derive explicit reconstruction formulas of the back-
projection type (see Theorem 1.1) as well as a Fourier slice identity (see
Theorem 1.3) similar to the one of the classical Radon transform.
### 1.1 Statement of the main results
Before we present our main results we introduce some notation. By
$C^{\infty}_{c}(\mathbb{R}^{d-1}\times\left(0,\infty\right))$ we denote the
space of all functions defined on $\mathbb{R}^{d}$, that are $C^{\infty}$ and
have compact support in $\mathbb{R}^{d-1}\times\left(0,\infty\right)$.
Likewise $C^{\infty}(\mathbb{R}^{d-1}\times\left(0,\pi/2\right))$ denotes the
space of all infinitely smooth functions defined on
$\mathbb{R}^{d-1}\times\left(0,\pi/2\right)$. As can easily be seen, the
conical Radon transform defined by (1.1) is well defined as an operator
$\mathcal{R}^{(p)}\colon
C_{c}^{\infty}\bigl{(}\mathbb{R}^{d-1}\times(0,\infty)\bigr{)}\to
C^{\infty}\bigl{(}\mathbb{R}^{d-1}\times(0,\pi/2)\bigr{)}$.
Points in $\mathbb{R}^{d}$ will be written in the form
$\left(\mathbf{x},y\right)$ with $\mathbf{x}\in\mathbb{R}^{d-1}$ and
$y\in\mathbb{R}$. The Fourier transform of a function
$f\colon\mathbb{R}^{d-1}\times\mathbb{R}\to\mathbb{C}$ with respect to the
first component is denoted by
$\left(\mathcal{F}f\right)\left(\mathbf{k},y\right)\coloneqq\int_{\mathbb{R}^{d-1}}e^{-i\mathbf{k}\mathbf{x}}f\left(\mathbf{x},y\right)\mathrm{d}\mathbf{x}$
for $\left(\mathbf{k},y\right)\in\mathbb{R}^{d-1}\times\mathbb{R}$. The Hankel
transform of order $\left(d-3\right)/2$ in the second argument is denoted by
$\left(\mathcal{H}_{(d-3)/2}f\right)\left(\mathbf{x},\lambda\right)\coloneqq\int_{0}^{\infty}J_{(d-3)/2}\left(y\lambda\right)f\left(\mathbf{x},y\right)y\mathrm{d}y$
for
$\left(\mathbf{x},\lambda\right)\in\mathbb{R}^{d-1}\times\left(0,\infty\right)$,
where $J_{(d-3)/2}$ is the Bessel function of the first kind of order
$\left(d-3\right)/2$. Note that for $d=2$, we have
$J_{-1/2}\left(y\right)=\sqrt{2/\pi y}\cos\left(y\right)$ and hence
$\mathcal{H}_{-1/2}$ is closely related to the cosine transform.
Similarly, we denote by $\mathcal{F}g$ the Fourier transform of a function
$g\colon\mathbb{R}^{d-1}\times\left(0,\pi/2\right)\to\mathbb{C}$ with respect
to the first argument. Finally, we denote by
$\mathcal{I}^{(1-d)}g\colon\mathbb{R}^{d-1}\times\left(0,\pi/2\right)\to\mathbb{C}$
the Riesz potential of $g$, defined by
$\left(\mathcal{F}\mathcal{I}^{(1-d)}g\right)\left(\mathbf{k},\theta\right)\coloneqq\left|\mathbf{k}\right|^{d-1}\left(\mathcal{F}g\right)\left(\mathbf{k},\theta\right)\quad\text{
for
}\left(\mathbf{k},\theta\right)\in\mathbb{R}^{d-1}\times\left(0,\frac{\pi}{2}\right)\,.$
(1.2)
The Riesz potential is well defined if
$\left(\mathcal{F}g\right)\left(\,\cdot\,,\theta\right)\in
L^{1}\bigl{(}\mathbb{R}^{d-1}\bigr{)}$ for every
$\theta\in\left(0,\pi/2\right)$, which will always be the case in our
considerations.
#### Explicit reconstruction formulas
The central results of this paper are the following explicit reconstruction
formulas for inverting the conical Radon transform.
###### Theorem 1.1 (Reconstruction formulas for the conical Radon transform).
For every $p\in\mathbb{R}$, every $f\in
C^{\infty}_{c}\bigl{(}\mathbb{R}^{d-1}\times\left(0,\infty\right)\bigr{)}$ and
every
$\left(\mathbf{x},y\right)\in\mathbb{R}^{d-1}\times\left(0,\infty\right)$, we
have
$\displaystyle f\left(\mathbf{x},y\right)$
$\displaystyle=\frac{y^{p}}{\left(2\pi\right)^{d-1}}\int_{0}^{\pi/2}\frac{1}{\cos\left(\theta\right)^{1+p}}\int_{S^{d-2}}\left(\mathcal{I}^{(1-d)}\mathcal{R}^{(p)}f\right)\left(\mathbf{x}+y\tan\left(\theta\right)\mathbf{n},\theta\right)\mathrm{d}\mathbf{n}\mathrm{d}\theta\,,$
(1.3) $\displaystyle f\left(\mathbf{x},y\right)$
$\displaystyle=\frac{1}{\left(2\pi\right)^{d-1}}\int_{\mathbb{R}^{d-1}}\frac{\left(\left|\mathbf{u}-\mathbf{x}\right|^{2}+y^{2}\right)^{\frac{p-1}{2}}}{\left|\mathbf{u}-\mathbf{x}\right|^{d-2}}\left(\mathcal{I}^{(1-d)}\mathcal{R}^{(p)}f\right)\left(\mathbf{u},\arctan\left(\frac{\left|\mathbf{u}-\mathbf{x}\right|}{y}\right)\right)\,\mathrm{d}\mathbf{u}\,.$
(1.4)
Here $\mathcal{I}^{(1-d)}$ is the Riesz potential defined by (1.2).
###### Proof.
See Sections 2.2 and 2.3. ∎
The reconstruction formulas (1.3), (1.4) are of the filtered back-projection
type: The Riesz potential can be interpreted as a filtration step in the first
argument and the integrations actually sum over all conical surfaces that pass
through the reconstruction point
$\left(\mathbf{x},y\right)\in\mathbb{R}^{d-1}\times\left(0,\infty\right)$. In
analogy to the classical Radon transform the integration process may therefore
be called _conical back-projection_. Note that (1.3), (1.4) only differ up to
a different parametrization of the set of all conical surfaces passing through
the reconstruction point.
For practical applications, the two and three dimensional situations are the
most relevant ones. In these cases the formulas of Theorem 1.1 read as
follows.
###### Corollary 1.2 (Reconstruction formulas for $d=2,3$).
1. (a)
Suppose $d=2$. Then, for every $f\in
C^{\infty}_{c}(\mathbb{R}\times\left(0,\infty\right))$ and every
$\left(x,y\right)\in\mathbb{R}\times\left(0,\infty\right)$,
$\displaystyle f\left(x,y\right)$
$\displaystyle=\frac{y^{p}}{2\pi}\int_{0}^{\pi/2}\frac{\bigl{(}\partial_{u}\mathrm{H}_{u}\mathcal{R}^{(p)}f\bigr{)}\left(x+\mathbf{n}y\tan\left(\theta\right),\theta\right)+\bigl{(}\partial_{u}\mathrm{H}_{u}\mathcal{R}^{(p)}f\bigr{)}\left(x-\mathbf{n}y\tan\left(\theta\right),\theta\right)}{\cos\left(\theta\right)^{1+p}}\,\mathrm{d}\theta\,,$
$\displaystyle f\left(x,y\right)$
$\displaystyle=\frac{1}{2\pi}\int_{\mathbb{R}}\left(\left|u-x\right|^{2}+y^{2}\right)^{\frac{p-1}{2}}\left(\partial_{u}\mathrm{H}_{u}\mathcal{R}^{(p)}f\right)\left(x,\arctan\left(\frac{\left|u-x\right|}{y}\right)\right)\,\mathrm{d}u\,.$
Here $\partial_{u}$ and $\mathrm{H}_{u}$ denote the derivative and the Hilbert
transform in first argument.
2. (b)
Suppose $d=3$. Then, for every $f\in
C^{\infty}_{c}\bigl{(}\mathbb{R}^{2}\times\left(0,\infty\right)\bigr{)}$ and
every $\left(\mathbf{x},y\right)\in\mathbb{R}^{2}\times\left(0,\infty\right)$,
$\displaystyle f\left(\mathbf{x},y\right)$
$\displaystyle=-\frac{y^{p}}{4\pi^{2}}\int_{0}^{\pi/2}\frac{1}{\cos\left(\theta\right)^{1+p}}\int_{S^{1}}\left(\Delta_{\mathbf{u}}\mathcal{R}^{(p)}f\right)\left(\mathbf{x}+y\tan\left(\theta\right)\mathbf{n},\theta\right)\mathrm{d}\mathbf{n}\mathrm{d}\theta$
$\displaystyle f\left(\mathbf{x},y\right)$
$\displaystyle=-\frac{1}{4\pi^{2}}\int_{\mathbb{R}^{2}}\frac{\left(\left|\mathbf{u}-\mathbf{x}\right|^{2}+y^{2}\right)^{\frac{p-1}{2}}}{\left|\mathbf{u}-\mathbf{x}\right|}\left(\Delta_{\mathbf{u}}\mathcal{R}^{(p)}f\right)\left(\mathbf{u},\arctan\left(\frac{\left|\mathbf{u}-\mathbf{x}\right|}{y}\right)\right)\,\mathrm{d}\mathbf{u}\,.$
Here $\Delta_{\mathbf{u}}$ denotes the Laplacian in the first component.
###### Proof.
For $d=2$ we have the Fourier representation
$\bigl{(}\mathcal{F}\mathrm{H}_{u}f\bigr{)}\left(k\right)=-i\operatorname{sign}\left(k\right)\bigl{(}\mathcal{F}f\bigr{)}\left(k\right)$
and
$\bigl{(}\mathcal{F}\partial_{u}f\bigr{)}\left(k\right)=ik\bigl{(}\mathcal{F}f\bigr{)}\left(k\right)$
of the Hilbert transform and the one dimensional derivative, respectively.
This shows $\mathcal{I}^{-1}=\partial_{u}\mathrm{H}_{u}$. Hence Item (a)
follows from Theorem 1.1. Similarly, for $d=3$, we have
$\mathcal{I}^{-2}=-\Delta_{\mathbf{u}}$ and hence Item (b) again follows from
Theorem 1.1. ∎
For $d=2$, formulas equivalent to the ones of Theorem 1.2 (a) have been first
derived in [BZG97, TN11]. The three dimensional reconstruction formulas of
Theorem 1.2 (b) (as well as the higher dimensional generalizations of Theorem
1.1) are new. One notes, that in three spatial dimensions the reconstruction
formulas are particularly simple and further local: The reconstruction of $f$
at some reconstruction point $\left(\mathbf{x},y\right)$ only requires the
integrals over cones passing through an arbitrarily small neighbourhood of
$\left(\mathbf{x},y\right)$. Since for any odd $d$, the Riesz potential
satisfies
$\mathcal{I}^{(1-d)}=\left(-1\right)^{(d-1)/2}\Delta_{\mathbf{u}}^{(d-1)/2}$,
the reconstruction formulas (1.3), (1.4) are in fact local for every odd space
dimensions. Contrary, in even space dimension (1.3), (1.4) are non-local:
Recovering a function at a single point requires knowledge of the integrals
over all conical surfaces. This behaviour is similar to the one of the
classical Radon transform, where also the inversion is local in odd and non-
local in even dimensions (see, for example, [Nat01, p. 20]).
#### A Fourier slice identity
Theorem 1.1 will be established using the following Theorem 1.3, which an
analogon of the well known Fourier slice identity [Nat01, Chapter 1, Theorem
1.1] satisfied by the classical Radon transform.
###### Theorem 1.3 (Fourier slice identity for the conical Radon transform).
For every $p\in\mathbb{R}$, every $f\in
C_{c}^{\infty}\bigl{(}\mathbb{R}^{d-1}\times(0,\infty)\bigr{)}$ and every
$\left(\mathbf{k},\theta\right)\in\mathbb{R}\times(0,\pi/2)$, we have
$\left(\mathcal{H}_{\frac{d-3}{2}}\mathcal{F}y^{\frac{d-3}{2}-p}f\right)\left(\mathbf{k},\left|\mathbf{k}\right|\tan\left(\theta\right)\right)=\left(2\pi\right)^{\frac{1-d}{2}}\frac{\cos\left(\theta\right)^{1-p}}{\tan\left(\theta\right)^{\frac{d-1}{2}}}\left|\mathbf{k}\right|^{\frac{d-3}{2}}\left(\mathcal{F}\mathcal{R}^{(p)}f\right)\left(\mathbf{k},\theta\right)\,.$
(1.5)
Here $y^{(d-3)/2-p}f$ is the function $\left(\mathbf{x},y\right)\mapsto
y^{(d-3)/2-p}f\left(\mathbf{x},y\right)$, $\mathcal{F}$ the Fourier transform
in the first argument, and $\mathcal{H}_{(d-3)/2}$ the Hankel transform of
order $(d-3)/2$ in the second argument.
###### Proof.
See Section 2.1. ∎
The Fourier slice identity is of course of interest on its own. The argument
$(\mathbf{k},\left|\mathbf{k}\right|\tan\theta)$, for
$\mathbf{k}\in\mathbb{R}^{d-1}$ and $\alpha\in(0,\pi/2)$, appearing on the
left hand side of (1.5), fills in the whole upper half-space, which is
required to invert the Fourier-Hankel transform using well known explicit and
stable inversion formulas. Hence the function $f$ can be reconstructed based
on (1.5) by means of a $d-1$-dimensional Fourier transform, followed by an
interpolation, and finally performing an inverse $d$-dimensional Fourier-
Hankel transform.
### 1.2 Outline
The remainder of the paper is mainly devoted to the proofs of Theorems 1.1 and
1.3 that we will establish in the following Section 2. We will first derive
the Fourier slice identity of Theorem 1.3, which will then be used to proof
the reconstruction formulas of Theorem 1.1. The paper ends with a discussion
in Section 3.
## 2 Proofs of the main results
In this section we derive Theorems 1.1 and 1.3. The following elementary Lemma
shows that it suffices to derive these results for the special case $p=0$.
###### Lemma 2.1 (Relation between $\mathcal{R}^{(p)}$ and
$\mathcal{R}^{(0)}$).
For every $p\in\mathbb{R}$, every $f\in
C_{c}^{\infty}\bigl{(}\mathbb{R}^{d-1}\times(0,\infty)\bigr{)}$ and every
$\left(\mathbf{k},\theta\right)\in\mathbb{R}^{d-1}\times(0,\pi/2)$, we have
$\left(\mathcal{R}^{(p)}f\right)\left(\mathbf{u},\theta\right)=\left(\cos\left(\theta\right)^{p}\mathcal{R}^{(0)}y^{-p}f\right)\left(\mathbf{u},\theta\right)\,.$
(2.1)
Here $y^{-p}$ stands for the operator that multiplies a function
$f\left(\mathbf{x},y\right)$ by $y^{-p}$ and likewise
$\cos\left(\theta\right)^{p}$ stands for the operator that multiplies
$g\left(\mathbf{u},\theta\right)$ by $\cos\left(\theta\right)^{p}$.
###### Proof.
The definition of $\mathcal{R}^{(p)}$ and the substitution
$s=y/\cos\left(\theta\right)$ yield
$\left(\mathcal{R}^{(p)}f\right)\left(\mathbf{u},\theta\right)=\int_{0}^{\infty}\frac{1}{s^{p}}\int_{S^{d-2}}f\left(\mathbf{u}+s\sin\left(\theta\right)\mathbf{n},s\cos\left(\theta\right)\right)\left(s\sin\left(\theta\right)\right)^{d-2}\mathrm{d}\mathbf{n}\mathrm{d}s\\\
=\cos\left(\theta\right)^{p-1}\int_{0}^{\infty}\int_{S^{d-2}}y^{-p}f\left(\mathbf{u}+y\tan\left(\theta\right)\mathbf{n},y\right)\left(y\tan\left(\theta\right)\right)^{d-2}\mathrm{d}\mathbf{n}\mathrm{d}y\,.$
Comparing the last expression for $\mathcal{R}^{(p)}f$ with the corresponding
expression for $\mathcal{R}^{(0)}f$ obviously shows (2.1). ∎
### 2.1 Proof of Theorem 1.3 (the Fourier slice identity)
We start by showing (1.5) for the special case $p=0$. The general case will
then be a consequence of Lemma 2.1.
The definition of the conical Radon transform, the definition of the Fourier
transform and some basic manipulations yield
$\displaystyle\left(\mathcal{F}\mathcal{R}^{(0)}f\right)\left(\mathbf{k},\theta\right)$
$\displaystyle=\int_{\mathbb{R}^{d-1}}e^{-i\mathbf{k}\mathbf{u}}\left(\mathcal{R}^{(0)}f\right)\left(\mathbf{u},\theta\right)\mathrm{d}\mathbf{u}$
$\displaystyle=\int_{0}^{\infty}\int_{S^{d-2}}\int_{\mathbb{R}^{d-1}}e^{-i\mathbf{k}\mathbf{u}}f\left(\mathbf{u}+s\sin\left(\theta\right)\mathbf{n},s\cos\left(\theta\right)\right)\left(s\sin\left(\theta\right)\right)^{d-2}\,\mathrm{d}\mathbf{u}\mathrm{d}\mathbf{n}\mathrm{d}s$
$\displaystyle=\int_{0}^{\infty}\left(s\sin\left(\theta\right)\right)^{d-2}\int_{S^{d-2}}e^{i\mathbf{k}s\sin\left(\theta\right)\mathbf{n}}\int_{\mathbb{R}^{d-1}}e^{-i\mathbf{k}\mathbf{u}}f\left(\mathbf{x},s\cos\left(\theta\right)\right)\,\mathrm{d}\mathbf{u}\mathrm{d}\mathbf{n}\mathrm{d}s$
$\displaystyle=\int_{0}^{\infty}\left(s\sin\left(\theta\right)\right)^{d-2}\left(\mathcal{F}f\right)\left(\mathbf{k},s\cos\left(\theta\right)\right)\left[\int_{S^{d-2}}e^{i\mathbf{k}s\sin\left(\theta\right)\mathbf{n}}\mathrm{d}\mathbf{n}\right]\mathrm{d}s\,.$
Now we use the identity (see, for example, [Nat01, page 198]),
$\int_{S^{d-2}}e^{-i\mathbf{k}r\mathbf{n}}\mathrm{d}\mathbf{n}=\left(2\pi\right)^{\frac{d-1}{2}}\left|\mathbf{k}\right|^{\frac{3-d}{2}}r^{\frac{3-d}{2}}J_{\frac{d-3}{2}}\left(\left|\mathbf{k}\right|r\right)\quad\text{
for all
}\left(\mathbf{k},r\right)\in\mathbb{R}^{d-1}\times\left(0,\infty\right)\,.$
(2.2)
Application of (2.2) with $r=s\sin\left(\theta\right)$ followed by the
substitution $s=y/\cos\left(\theta\right)$ yields
$\displaystyle\left(\mathcal{F}\mathcal{R}^{(0)}f\right)\left(\mathbf{k},\theta\right)$
$\displaystyle=\left(2\pi\right)^{\frac{d-1}{2}}\int_{0}^{\infty}\left(s\sin\left(\theta\right)\right)^{d-2}\left(\mathcal{F}f\right)\left(\mathbf{k},s\cos\left(\theta\right)\right)\left|\mathbf{k}\right|^{\frac{3-d}{2}}\left(s\sin\left(\theta\right)\right)^{\frac{3-d}{2}}J_{\frac{d-3}{2}}\left(\left|\mathbf{k}\right|s\sin\left(\theta\right)\right)\mathrm{d}s$
$\displaystyle=\left(2\pi\right)^{\frac{d-1}{2}}\int_{0}^{\infty}\left(s\sin\left(\theta\right)\right)^{\frac{d-1}{2}}\left(\mathcal{F}f\right)\left(\mathbf{k},s\cos\left(\theta\right)\right)\left|\mathbf{k}\right|^{\frac{3-d}{2}}J_{\frac{d-3}{2}}\left(\left|\mathbf{k}\right|s\sin\left(\theta\right)\right)\mathrm{d}s$
$\displaystyle=\left(2\pi\right)^{\frac{d-1}{2}}\int_{0}^{\infty}\left(y\tan\left(\theta\right)\right)^{\frac{d-1}{2}}\left(\mathcal{F}f\right)\left(\mathbf{k},y\right)\left|\mathbf{k}\right|^{\frac{3-d}{2}}J_{\frac{d-3}{2}}\left(\left|\mathbf{k}\right|y\tan\left(\theta\right)\right)\frac{\mathrm{d}y}{\cos\left(\theta\right)}$
$\displaystyle=\left(2\pi\right)^{\frac{d-1}{2}}\frac{\tan\left(\theta\right)^{\frac{d-1}{2}}}{\cos\left(\theta\right)}\left|\mathbf{k}\right|^{\frac{3-d}{2}}\int_{0}^{\infty}y^{\frac{d-3}{2}}\left(\mathcal{F}f\right)\left(\mathbf{k},y\right)J_{\frac{d-3}{2}}\left(\left|\mathbf{k}\right|y\tan\left(\theta\right)\right)y\mathrm{d}y\,.$
The last displayed equation is recognised as the Hankel transform of order
$\left(d-3\right)/2$ of $\mathcal{F}f$ in the second argument. We conclude,
that
$\left(\mathcal{F}\mathcal{R}^{(0)}f\right)\left(\mathbf{k},\theta\right)=\left(2\pi\right)^{\frac{d-1}{2}}\frac{\tan\left(\theta\right)^{\frac{d-1}{2}}}{\cos\left(\theta\right)}\left|\mathbf{k}\right|^{\frac{3-d}{2}}\left(\mathcal{H}_{\frac{d-3}{2}}\mathcal{F}y^{\frac{d-3}{2}}f\right)\left(\mathbf{k},\left|\mathbf{k}\right|\tan\left(\theta\right)\right)\,.$
(2.3)
This shows (1.5) for the special case $p=0$.
For general $p\in\mathbb{R}$ we use the relation
$\mathcal{R}^{(0)}y^{-p}f=\cos\left(\theta\right)^{-p}\mathcal{R}^{(p)}f$ from
Lemma 2.1. Together with (2.3) this yields
$\left(\mathcal{H}_{\frac{d-3}{2}}\mathcal{F}y^{\frac{d-3}{2}-p}f\right)\left(\mathbf{k},\left|\mathbf{k}\right|\tan\left(\theta\right)\right)=\left(\mathcal{H}_{\frac{d-3}{2}}\mathcal{F}y^{\frac{d-3}{2}}y^{-p}f\right)\left(\mathbf{k},\left|\mathbf{k}\right|\tan\left(\theta\right)\right)\\\
=\left(2\pi\right)^{\frac{1-d}{2}}\frac{\cos\left(\theta\right)}{\tan\left(\theta\right)^{\frac{d-1}{2}}}\left|\mathbf{k}\right|^{\frac{d-3}{2}}\left(\mathcal{F}\mathcal{R}^{(0)}y^{-p}f\right)\left(\mathbf{k},\theta\right)=\left(2\pi\right)^{\frac{1-d}{2}}\frac{\cos\left(\theta\right)^{1-p}}{\tan\left(\theta\right)^{\frac{d-1}{2}}}\left|\mathbf{k}\right|^{\frac{d-3}{2}}\left(\mathcal{F}\mathcal{R}^{(p)}f\right)\left(\mathbf{k},\theta\right)\,.$
This is (1.5) for the case of general $p\in\mathbb{R}$ and concludes the proof
of Theorem 1.3.
### 2.2 Proof of reconstruction formula (1.3)
We start with the proof of (1.3) for $p=0$. Application of the inversion
formulas for the Fourier and the Hankel transform followed by the substitution
$\lambda=\left|\mathbf{k}\right|\tan\left(\theta\right)$ shows
$\displaystyle y^{\frac{d-3}{2}}f\left(\mathbf{x},y\right)$
$\displaystyle=\frac{1}{\left(2\pi\right)^{d-1}}\int_{\mathbb{R}^{d-1}}\int_{0}^{\infty}\left(\mathcal{H}_{\frac{d-3}{2}}\mathcal{F}y^{\frac{d-3}{2}}f\right)\left(\mathbf{k},\lambda\right)J_{\frac{d-3}{2}}\left(\lambda
y\right)e^{i\mathbf{k}\mathbf{x}}\,\lambda\mathrm{d}\lambda\,\mathrm{d}\mathbf{k}$
$\displaystyle=\frac{1}{\left(2\pi\right)^{d-1}}\int_{\mathbb{R}^{d-1}}\int_{0}^{\pi/2}\left(\mathcal{H}_{\frac{d-3}{2}}\mathcal{F}y^{\frac{d-3}{2}}f\right)\left(\mathbf{k},\tan\left(\theta\right)y\right)J_{\frac{d-3}{2}}\left(\left|\mathbf{k}\right|\tan\left(\theta\right)y\right)e^{i\mathbf{k}\mathbf{x}}\frac{\left|\mathbf{k}\right|^{2}\tan\left(\theta\right)}{\cos\left(\theta\right)^{2}}\,\mathrm{d}\theta\mathrm{d}\mathbf{k}\,.$
Application of the Fourier slice identity (Theorem 1.3) with $p=0$ and
interchanging the order of integration then yields
$y^{\frac{d-3}{2}}f\left(\mathbf{x},y\right)=\left(2\pi\right)^{\frac{3\left(1-d\right)}{2}}\int_{0}^{\pi/2}\frac{\left(\tan\left(\theta\right)\right)^{\frac{3-d}{2}}}{\cos\left(\theta\right)}\int_{\mathbb{R}^{d-1}}\left|\mathbf{k}\right|^{\frac{d+1}{2}}\left(\mathcal{F}\mathcal{R}^{(0)}f\right)\left(\mathbf{k},\theta\right)J_{\frac{d-3}{2}}\left(\left|\mathbf{k}\right|\tan\left(\theta\right)y\right)e^{i\mathbf{k}\mathbf{x}}\mathrm{d}\mathbf{k}\,\mathrm{d}\theta\,.$
(2.4)
By (2.2), we have
$J_{\frac{d-3}{2}}\left(\left|\mathbf{k}\right|\tan\left(\theta\right)y\right)=\left(2\pi\right)^{\frac{1-d}{2}}\left|\mathbf{k}\right|^{\frac{d-3}{2}}\left(\tan\left(\theta\right)y\right)^{\frac{d-3}{2}}\int_{S^{d-2}}e^{-i\mathbf{k}\tan\left(\theta\right)y\mathbf{n}}\mathrm{d}\mathbf{n}\,.$
Therefore,
$\int_{\mathbb{R}^{d-1}}\left|\mathbf{k}\right|^{\frac{d+1}{2}}\left(\mathcal{F}\mathcal{R}^{(0)}f\right)\left(\mathbf{k},\theta\right)J_{\frac{d-3}{2}}\left(\left|\mathbf{k}\right|\tan\left(\theta\right)y\right)e^{i\mathbf{k}\mathbf{x}}\mathrm{d}\mathbf{k}\\\
\begin{aligned}
&=\left(2\pi\right)^{\frac{1-d}{2}}\int_{S^{d-2}}\left(\tan\left(\theta\right)y\right)^{\frac{d-3}{2}}\left[\int_{\mathbb{R}^{d-1}}\left|\mathbf{k}\right|^{d-1}\left(\mathcal{F}\mathcal{R}^{(0)}f\right)\left(\mathbf{k},\theta\right)e^{i\mathbf{k}\left(\mathbf{x}-\tan\left(\theta\right)y\mathbf{n}\right)}\mathrm{d}\mathbf{k}\right]\\\
&=\left(2\pi\right)^{\frac{d-1}{2}}\int_{S^{d-2}}\left(\tan\left(\theta\right)y\right)^{\frac{d-3}{2}}\left(\mathcal{I}^{(1-d)}\mathcal{R}^{(0)}f\right)\left(\mathbf{x}-\tan\left(\theta\right)y\mathbf{n},\theta\right)\mathrm{d}\mathbf{n}\,.\end{aligned}$
Together with (2.4) this further implies
$y^{\frac{d-3}{2}}f\left(\mathbf{x},y\right)=\frac{1}{\left(2\pi\right)^{d-1}}\int_{0}^{\pi/2}\frac{\left(\tan\left(\theta\right)\right)^{\frac{3-d}{2}}}{\cos\left(\theta\right)}\int_{S^{d-2}}\left(\tan\left(\theta\right)y\right)^{\frac{d-3}{2}}\left(\mathcal{I}^{(1-d)}\mathcal{R}^{(0)}f\right)\left(\mathbf{x}-\tan\left(\theta\right)y\mathbf{n},\theta\right)\mathrm{d}\mathbf{n}\,\mathrm{d}\theta\\\
=\frac{y^{\frac{d-3}{2}}}{\left(2\pi\right)^{d-1}}\int_{0}^{\pi/2}\frac{1}{\cos\left(\theta\right)}\int_{S^{d-2}}\left(\mathcal{I}^{(1-d)}\mathcal{R}^{(0)}f\right)\left(\mathbf{x}-\tan\left(\theta\right)y\mathbf{n},\theta\right)\mathrm{d}\mathbf{n}\,\mathrm{d}\theta\,.$
This shows formula (1.3) for the special case $p=0$.
To show (1.3) in the general case $p\in\mathbb{R}$, we again use the relation
$\cos\left(\theta\right)^{-p}\mathcal{R}^{(p)}f=\mathcal{R}^{(0)}y^{-p}f$ from
Lemma 2.1. Hence application of the reconstruction formula for the special
case $p=0$ to $\mathcal{R}^{(0)}y^{-p}f$ yields
$y^{-p}f\left(\mathbf{x},y\right)=\left(2\pi\right)^{1-d}\int_{0}^{\pi/2}\frac{1}{\cos\left(\theta\right)^{1+p}}\int_{S^{d-2}}\left(\mathcal{I}^{(1-d)}\mathcal{R}^{(p)}f\right)\left(\mathbf{x}-\tan\left(\theta\right)y\mathbf{n},\theta\right)\mathrm{d}\mathbf{n}\,\mathrm{d}\theta\,.$
This shows (1.3) in the general case $\mu\in\mathbb{R}$.
### 2.3 Proof of reconstruction formula (1.4)
Finally we derive (1.4) as an easy consequence of (1.3). To that end we first
substitute $\theta=\arctan\left(r/y\right)$ with $r\in\left(0,\infty\right)$.
Then $\mathrm{d}\theta=y^{-1}\cos\left(\theta\right)^{2}\mathrm{d}r$ and
$\cos\left(\theta\right)=1/\sqrt{1+r^{2}/y^{2}}$. Consequently, (1.3) implies
$f\left(\mathbf{x},y\right)=\frac{y^{p-1}}{\left(2\pi\right)^{d-1}}\int_{0}^{\infty}\int_{S^{d-2}}\left(1+r^{2}/y^{2}\right)^{\frac{p-1}{2}}\left(\mathcal{I}^{(1-d)}\mathcal{R}^{(0)}f\right)\left(\mathbf{x}-r\mathbf{n},\arctan\left(\frac{r}{y}\right)\right)\mathrm{d}\mathbf{n}\,\mathrm{d}r\\\
=\frac{1}{\left(2\pi\right)^{d-1}}\int_{0}^{\infty}\int_{S^{d-2}}\left(r^{2}+y^{2}\right)^{\frac{p-1}{2}}\left(\mathcal{I}^{(1-d)}\mathcal{R}^{(0)}f\right)\left(\mathbf{x}+r\mathbf{n},\arctan\left(\frac{r}{y}\right)\right)\mathrm{d}\mathbf{n}\,\mathrm{d}r\,.$
Now we substitute $\mathbf{x}+r\mathbf{n}=\mathbf{u}$ (polar coordinates in
the plane $\mathbb{R}^{d-1}$ around the center $\mathbf{x}$). Then
$\mathrm{d}\mathbf{u}=r^{d-2}\mathrm{d}\mathbf{n}\mathrm{d}r$ and
$r=\left|\mathbf{u}-\mathbf{x}\right|$. Consequently,
$f\left(\mathbf{x},y\right)=\frac{1}{\left(2\pi\right)^{d-1}}\int_{\mathbb{R}^{d-1}}\frac{\left(\left|\mathbf{u}-\mathbf{x}\right|^{2}+y^{2}\right)^{\frac{p-1}{2}}}{\left|\mathbf{u}-\mathbf{x}\right|^{d-2}}\left(\mathcal{I}^{(1-d)}\mathcal{R}^{(0)}f\right)\left(\mathbf{x},\arctan\left(\frac{\left|\mathbf{u}-\mathbf{x}\right|}{y}\right)\right)\;\mathrm{d}\mathbf{x}\,.$
This is the reconstruction formula (1.4).
## 3 Discussion
In this paper we derived explicit reconstruction formulas for the conical
Radon transform, which integrates a function in $d$ spatial variables over all
cones with vertices on a hyperplane and symmetry axis orthogonal to this
plane. The derived formulas are of the back-projection type and are
theoretically exact. Further, they are local for odd $d$, and non-local for
even $d$. Among others, inversion of the conical Radon transform is relevant
for emission tomography using Compton cameras as proposed in [EFTN77, Sin83,
TNE74]. Such a device measures the direction as well as the scattering angle
of an incoming photon at the front of the camera. The location of the photon
emission can therefore be traced back to the surface of a cone. Recovering the
density of the photon source therefore yields to the inversion of the conical
Radon transform in a natural manner.
Radon transforms are the theoretical foundation of many medical imaging and
remote sensing application. Certainly the most well known instance is the
classical Radon transform, which integrates a function over hyperplanes. Among
others, inversion of the classical Radon transform is important for classical
transmission computed tomography and has been studied in many textbooks (see,
for example, [Hel99, Nat01]). Closed form reconstruction formulas are known
for a long time and have first been derived already in 1917 by J. Radon
[Rad17]. Another Radon transform that has been studied in detail more recently
is the spherical Radon transform. This transform integrates a function over
spherical surfaces (for some restricted centers of integration) and is, among
others, important for photo- and thermoacoustic tomography [KK08]. Closed form
reconstruction formulas for planar and spherical center sets have been found
in [And88, Kun07, Faw85, FHR07, FPR04]. The conical Radon transform, on the
other hand, is much less studied. In particular, closed form reconstruction
formulas have only been known for the case $d=2$, see [BZG97, TN11]. In this
paper we derived such reconstruction formulas for arbitrary dimension $d\geq
2$. For computed tomography with Compton cameras [BZG98, CB94], the three
dimensional case is of course the most relevant one. In this case, our
reconstruction formulas have a particularly simple structure and consist of an
application of the Laplacian followed by a conical back-projection. The
numerical implementation seems quite straight forward following the ones of
the classical or the spherical Radon transform (see, for example, [FHR07,
Nat01]). Numerical studies, however, will be subject of future research.
## References
* [And88] L. E. Andersson. On the determination of a function from spherical averages. SIAM J. Math. Anal., 19(1):214–232, 1988.
* [BZG97] R. Basko, G. L. Zeng, and G. T. Gullberg. Analytical reconstruction formula for one-dimensional compton camera. IEEE Trans. Nucl. Sci., 44(3):1342–1346, 1997.
* [BZG98] R. Basko, G. L. Zeng, and G. T. Gullberg. Application of spherical harmonics to image reconstruction for the compton camera. Physics in Medicine and Biology, 43(4):887, 1998.
* [CB94] M. J. Cree and P. J. Bones. Towards direct reconstruction from a gamma camera based on compton scattering. IEEE Trans. Med. Imag., 13(2):398–407, 1994.
* [EFTN77] D. B. Everett, J. S. Fleming, R. W. Todd, and J. M. Nightingale. Gamma-radiation imaging system based on the compton effect. Proc. IEEE, 124(11):995–1000, 1977.
* [Faw85] J. A. Fawcett. Inversion of $n$-dimensional spherical averages. SIAM J. Appl. Math., 45(2):336–341, 1985.
* [FHR07] D. Finch, M. Haltmeier, and Rakesh. Inversion of spherical means and the wave equation in even dimensions. SIAM J. Appl. Math., 68(2):392–412, 2007.
* [FPR04] D. Finch, S. Patch, and Rakesh. Determining a function from its mean values over a family of spheres. SIAM J. Math. Anal., 35(5):1213–1240, 2004.
* [Hel99] S. Helgason. The Radon Transform, volume 5 of Progress in Mathematics. Birkhäuser, Boston, second edition, 1999.
* [KK08] P. Kuchment and L. A. Kunyansky. Mathematics of thermoacoustic and photoacoustic tomography. European J. Appl. Math., 19:191–224, 2008.
* [Kun07] L. A. Kunyansky. Explicit inversion formulae for the spherical mean Radon transform. Inverse Probl., 23(1):373–383, 2007.
* [Nat01] F. Natterer. The Mathematics of Computerized Tomography, volume 32 of Classics in Applied Mathematics. SIAM, Philadelphia, 2001.
* [NTG05] M. K. Nguyen, T. T. Truong, and P. Grangeat. Radon transforms on a class of cones with fixed axis direction. J. Phys. A, 38(37):8003–8015, 2005.
* [Par00] L. C. Parra. Reconstruction of cone-beam projections from compton scattered data. IEEE Trans. Nucl. Sci., 47(4):1543–1550, 2000.
* [Rad17] J. Radon. Über die Bestimmung von Funktionen durch ihre Integralwerte längs gewisser Mannigfaltigkeiten. Ber. Verh. Kön. Sächs. Ges. Wiss. Leipzig Math. Phys. Kl., 69:262–277, 1917.
* [Sin83] M. Singh. An electronically collimated gamma camera for single photon emission computed tomography. Part I: Theoretical considerations and design criteria. Med. Phys., 10(421):1983, 1983.
* [TN11] T. T. Truong and M. K. Nguyen. On new V-line Radon transforms in $\mathbb{R}^{2}$ and their inversion. J. Phys. A, 44(7):075206, 13, 2011.
* [TNE74] R. W. Todd, J. M. Nightingale, and D. B. Everett. A proposed gamma camera. Nature, 251:132–134, 1974.
|
arxiv-papers
| 2013-10-01T11:06:44 |
2024-09-04T02:49:51.787451
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Markus Haltmeier",
"submitter": "Markus Haltmeier",
"url": "https://arxiv.org/abs/1310.0239"
}
|
1310.0261
|
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k1,39:13552292,18785015:274760 x1,39:15153971,18785015
k1,39:15428730,18785015:274759 x1,39:17012065,18785015
k1,39:17547984,18785015:535919 x1,39:19112968,18785015
k1,39:19112968,18785015:0 ) (1,40:3729360,19571447:15383608,446948,141880
x1,40:4020336,19571447 x1,40:4338192,19571447 x1,40:4956848,19571447
x1,40:5357263,19571447 x1,40:6922893,19571447 k1,39:7182172,19571447:259279
x1,39:7800828,19571447 x1,39:8115404,19571447 x1,39:8396556,19571447
k1,39:8655835,19571447:259279 x1,39:9893147,19571447
k1,39:10152425,19571447:259278 x1,39:11426441,19571447
k1,39:11685720,19571447:259279 x1,39:12013400,19571447 x1,39:12331256,19571447
x1,39:12830640,19571447 k1,39:13089919,19571447:259279 x1,39:13380895,19571447
k1,39:13640174,19571447:259279 x1,39:14914181,19571447
k1,39:15173459,19571447:259278 x1,39:16629003,19571447
k1,39:16888282,19571447:259279 x1,39:17434194,19571447
k1,39:17693473,19571447:259279 x1,39:19112968,19571447
k1,39:19112968,19571447:0 ) (1,40:3729360,20357879:15383608,446948,141880
x1,40:4348016,20357879 x1,40:4638992,20357879 x1,40:4857224,20357879
x1,40:6065039,20357879 k1,39:6330027,20357879:264988 x1,39:7495232,20357879
x1,39:9296791,20357879 k1,39:9561778,20357879:264987 x1,39:9889458,20357879
x1,39:11150370,20357879 k1,39:11395129,20357879:244759 x1,39:12013785,20357879
x1,39:13669856,20357879 k1,39:13934844,20357879:264988 x1,39:14881180,20357879
k1,39:15125939,20357879:244759 x1,39:15854681,20357879
k1,39:16099439,20357879:244758 x1,39:16900278,20357879 x1,39:18338117,20357879
k1,39:18603105,20357879:264988 x1,39:19112968,20357879
k1,39:19112968,20357879:0 ) (1,40:3729360,21144311:15383608,446948,141880
x1,40:4057040,21144311 x1,39:5363814,21144311 k1,39:5501096,21144311:137282
x1,39:6083704,21144311 x1,39:6649928,21144311 x1,39:7923926,21144311
k1,39:8061208,21144311:137282 x1,39:8571071,21144311 x1,39:9116983,21144311
x1,39:10136054,21144311 x1,39:11155780,21144311 k1,39:11293062,21144311:137282
x1,39:11802925,21144311 k1,39:11940206,21144311:137281 x1,39:12558862,21144311
k1,39:12696144,21144311:137282 x1,39:13787959,21144311
k1,39:13925241,21144311:137282 x1,39:14435104,21144311
k1,39:14572386,21144311:137282 x1,39:15810344,21144311
k1,39:15947626,21144311:137282 x1,39:16857258,21144311 x1,39:17657442,21144311
x1,39:18130601,21144311 x1,39:19112968,21144311 k1,39:19112968,21144311:0 )
(1,40:3729360,21930743:15383608,446948,9168 x1,40:4530199,21930743
k1,39:4706579,21930743:176380 x1,39:6344306,21930743 x1,39:8673424,21930743
x1,39:10129623,21930743 k1,39:10306003,21930743:176380 x1,39:10815866,21930743
x1,39:11289025,21930743 x1,39:12089209,21930743 x1,39:13108935,21930743
k1,39:13285315,21930743:176380 x1,39:14231651,21930743
k1,39:14408030,21930743:176379 x1,39:14735710,21930743 x1,39:15192485,21930743
x1,39:15374668,21930743 x1,39:16466483,21930743 k1,39:16642863,21930743:176380
x1,39:17443702,21930743 k1,39:17620082,21930743:176380 x1,39:18567056,21930743
x1,39:18894736,21930743 x1,39:19112968,21930743 k1,40:19112968,21930743:0 )
(1,40:3729360,22717175:15383608,446948,141880 x1,40:4202519,22717175
x1,40:5222245,22717175 k1,39:5362079,22717175:139834 x1,39:6235671,22717175
x1,39:8001190,22717175 k1,39:8196349,22717175:195159 x1,39:9215429,22717175
k1,39:9355263,22717175:139834 x1,39:10810807,22717175
k1,39:10950642,22717175:139835 x1,39:11496554,22717175
k1,39:11636388,22717175:139834 x1,39:13055883,22717175
k1,39:13195718,22717175:139835 x1,39:14142054,22717175 x1,39:14651917,22717175
x1,39:15926571,22717175 k1,39:16066405,22717175:139834 x1,39:16503516,22717175
k1,39:16643351,22717175:139835 x1,39:17589687,22717175 x1,39:18463271,22717175
k1,39:18603105,22717175:139834 x1,39:19112968,22717175
k1,39:19112968,22717175:0 ) (1,40:3729360,23503607:15383608,446948,141880
x1,40:4348016,23503607 k1,39:4496627,23503607:148611 x1,39:5042531,23503607
k1,39:5191142,23503607:148611 x1,39:6537885,23503607
k1,39:6686496,23503607:148611 x1,39:7232400,23503607
k1,39:7381010,23503607:148610 x1,39:8291298,23503607 x1,39:8728409,23503607
x1,39:9529248,23503607 k1,39:9677859,23503607:148611 x1,39:10187722,23503607
k1,39:10336333,23503607:148611 x1,39:11719780,23503607
k1,39:11868391,23503607:148611 x1,39:12414303,23503607
k1,39:12562914,23503607:148611 x1,39:13363753,23503607
k1,39:13512364,23503607:148611 x1,39:14568140,23503607
k1,39:14716735,23503607:148595 x1,39:16172934,23503607 x1,39:17010477,23503607
k1,39:17159088,23503607:148611 x1,39:18378703,23503607
k1,39:18530360,23503607:151657 x1,39:19112968,23503607
k1,39:19112968,23503607:0 ) (1,40:3729360,24290039:15383608,446948,141880
x1,40:4712382,24290039 k1,39:4898634,24290039:186252 x1,39:5699473,24290039
k1,39:5885725,24290039:186252 x1,39:6213405,24290039 x1,39:6670180,24290039
x1,39:6852363,24290039 x1,39:8271858,24290039 k1,39:8458110,24290039:186252
x1,39:9405084,24290039 x1,39:9732764,24290039 x1,39:10205923,24290039
x1,39:11225649,24290039 k1,39:11411901,24290039:186252 x1,39:12285493,24290039
x1,39:13887172,24290039 k1,39:14073425,24290039:186253 x1,39:14983057,24290039
k1,39:15169309,24290039:186252 x1,39:15787965,24290039
k1,39:15974217,24290039:186252 x1,39:17175481,24290039
k1,39:17361733,24290039:186252 x1,39:17871596,24290039
k1,39:18057848,24290039:186252 x1,39:18894736,24290039 x1,39:19112968,24290039
k1,40:19112968,24290039:0 ) (1,40:3729360,25076471:15383608,446948,141880
x1,40:4712382,25076471 k1,39:4870161,25076471:157779 x1,39:5671000,25076471
k1,39:5828779,25076471:157779 x1,39:6447435,25076471 x1,39:7211579,25076471
x1,39:9140282,25076471 k1,39:9298061,25076471:157779 x1,39:9843973,25076471
k1,39:10001752,25076471:157779 x1,39:11348495,25076471
k1,39:11506275,25076471:157780 x1,39:12816978,25076471
k1,39:12974757,25076471:157779 x1,39:14139971,25076471
k1,39:14297750,25076471:157779 x1,39:14515982,25076471 x1,39:14806958,25076471
x1,39:15124814,25076471 x1,39:15958429,25076471 k1,39:16116208,25076471:157779
x1,39:16626071,25076471 k1,39:16783850,25076471:157779 x1,39:19112968,25076471
k1,40:19112968,25076471:0 ) (1,40:3729360,25862903:15383608,446948,141880
x1,40:5403791,25862903 x1,40:7186997,25862903 g1,39:7350837,25862903
x1,39:8297173,25862903 x1,39:9462396,25862903 x1,39:10535868,25862903
k1,40:19112968,25862903:8577100 g1,40:19112968,25862903 )
(1,42:3729360,26725001:15383608,446948,145636
h1,41:3729360,26725001:655361,0,0 x1,41:5950360,26725001
k1,41:6207253,26725001:256893 x1,41:7008092,26725001
k1,41:7264984,26725001:256892 x1,41:7883640,26725001 x1,41:8340424,26725001
k1,41:8597317,26725001:256893 x1,41:9544291,26725001 x1,41:9871971,26725001
x1,41:10345130,26725001 x1,41:11364856,26725001 k1,41:11621749,26725001:256893
x1,41:12859061,26725001 k1,41:13115953,26725001:256892 x1,41:15868441,26725001
k1,41:16125240,26725001:256799 x1,41:16635103,26725001
k1,41:16891996,26725001:256893 x1,41:17838970,26725001 x1,41:18166650,26725001
x1,41:19112968,26725001 k1,41:19112968,26725001:0 )
(1,42:3729360,27511433:15383608,446948,141880 x1,42:4530199,27511433
k1,41:4707327,27511433:177128 x1,41:5617615,27511433 x1,41:7837958,27511433
k1,41:8015086,27511433:177128 x1,41:8524294,27511433 x1,41:9034157,27511433
x1,41:9325133,27511433 x1,41:10344859,27511433 k1,41:10521988,27511433:177129
x1,41:10812964,27511433 x1,41:11140644,27511433 x1,41:11624298,27511433
x1,41:12497890,27511433 x1,41:13972432,27511433 k1,41:14152882,27511433:180450
x1,41:14443858,27511433 k1,41:14620986,27511433:177128 x1,41:14948666,27511433
x1,41:15514234,27511433 x1,41:16679457,27511433 k1,41:16856585,27511433:177128
x1,41:18276080,27511433 x1,41:19112968,27511433 k1,41:19112968,27511433:0 )
(1,42:3729360,28297865:15383608,446948,141880 x1,42:4348016,28297865
x1,42:5440487,28297865 k1,41:5587413,28297865:146926 x1,41:6460997,28297865
k1,41:6607923,28297865:146926 x1,41:7845235,28297865
k1,41:7992162,28297865:146927 x1,41:8319842,28297865 x1,41:8610818,28297865
x1,41:8922114,28297865 x1,41:9385449,28297865 x1,41:10659465,28297865
k1,41:10806391,28297865:146926 x1,41:11570535,28297865
k1,41:11717461,28297865:146926 x1,41:12518300,28297865
k1,41:12665226,28297865:146926 x1,41:13721002,28297865
k1,41:13867912,28297865:146910 x1,41:15324111,28297865 x1,41:16161654,28297865
k1,41:16308580,28297865:146926 x1,41:17364355,28297865
k1,41:17511281,28297865:146926 x1,41:19112968,28297865
k1,41:19112968,28297865:0 ) (1,42:3729360,29084297:15383608,446948,141880
x1,42:5112151,29084297 x1,41:5423447,29084297 x1,41:5959527,29084297
k1,41:6211703,29084297:252176 x1,41:7012542,29084297
k1,41:7264718,29084297:252176 x1,41:7773926,29084297 x1,41:9448357,29084297
k1,41:9700534,29084297:252177 x1,41:10173693,29084297 x1,41:10484989,29084297
x1,41:10948324,29084297 k1,41:11200500,29084297:252176 x1,41:11746412,29084297
k1,41:11998588,29084297:252176 x1,41:12981628,29084297 x1,41:13491491,29084297
x1,41:14546602,29084297 x1,41:15602377,29084297 k1,41:15854553,29084297:252176
x1,41:16764841,29084297 x1,41:17856657,29084297 x1,41:19112968,29084297
k1,41:19112968,29084297:0 ) (1,42:3729360,29870729:15383608,446948,141880
x1,42:4129784,29870729 x1,41:4405039,29870729 k1,41:4679532,29870729:274493
x1,41:5298188,29870729 x1,41:5771347,29870729 x1,41:7300281,29870729
k1,41:7574773,29870729:274492 x1,41:8375612,29870729
k1,41:8650105,29870729:274493 x1,41:8941081,29870729 x1,41:9258937,29870729
x1,41:10969408,29870729 k1,41:11243901,29870729:274493 x1,41:12299012,29870729
x1,41:13354787,29870729 k1,41:13629279,29870729:274492 x1,41:15085478,29870729
k1,41:15359971,29870729:274493 x1,41:16306307,29870729 x1,41:17253281,29870729
k1,41:17527773,29870729:274492 x1,41:18037636,29870729
k1,41:18312129,29870729:274493 x1,41:19112968,29870729
k1,41:19112968,29870729:0 ) (1,42:3729360,30657161:15383608,446948,141880
x1,42:4948975,30657161 k1,41:5135197,30657161:186222 x1,41:5972741,30657161
k1,41:6154469,30657161:181728 x1,41:6482149,30657161 x1,41:8144131,30657161
k1,41:8325876,30657161:181745 x1,41:9272850,30657161 x1,41:9600530,30657161
x1,41:10073689,30657161 x1,41:11348343,30657161 k1,41:11530089,30657161:181746
x1,41:12148745,30657161 x1,41:12463321,30657161 x1,41:12744473,30657161
k1,41:12926218,30657161:181745 x1,41:14163530,30657161
k1,41:14345276,30657161:181746 x1,41:15619292,30657161
k1,41:15801038,30657161:181746 x1,41:17293278,30657161
k1,41:17475023,30657161:181745 x1,41:18130383,30657161
k1,41:18312129,30657161:181746 x1,41:19112968,30657161
k1,41:19112968,30657161:0 ) (1,42:3729360,31443593:15383608,446948,141880
x1,42:4676334,31443593 x1,42:5004014,31443593 x1,42:6278012,31443593
k1,41:6485830,31443593:207818 x1,41:6995038,31443593 x1,41:7504901,31443593
x1,41:7795877,31443593 x1,41:8815603,31443593 k1,41:9023421,31443593:207818
x1,41:9314397,31443593 x1,41:9642077,31443593 x1,41:10125731,31443593
x1,41:10999323,31443593 x1,41:12310025,31443593 k1,41:12517843,31443593:207818
x1,41:13464179,31443593 k1,41:13671996,31443593:207817 x1,41:14327356,31443593
k1,41:14535174,31443593:207818 x1,41:15336013,31443593
k1,41:15543831,31443593:207818 x1,41:15834807,31443593 x1,41:16152663,31443593
x1,41:17863134,31443593 k1,41:18070952,31443593:207818 x1,41:18689608,31443593
x1,41:18907840,31443593 x1,41:19112968,31443593 k1,42:19112968,31443593:0 )
(1,42:3729360,32230025:15383608,446948,141880 x1,42:4275272,32230025
x1,42:6203975,32230025 k1,41:6348218,32230025:144243 x1,41:6894130,32230025
k1,41:7038373,32230025:144243 x1,41:7839212,32230025
k1,41:7983455,32230025:144243 x1,41:9039231,32230025
k1,41:9183455,32230025:144224 x1,41:10639654,32230025 x1,41:11477197,32230025
k1,41:11621440,32230025:144243 x1,41:12677215,32230025
k1,41:12821458,32230025:144243 x1,41:14240953,32230025 x1,41:15077841,32230025
k1,41:15222084,32230025:144243 x1,41:16096313,32230025 x1,41:16278496,32230025
x1,41:16606176,32230025 x1,41:16893879,32230025 x1,41:17913605,32230025
k1,41:18057848,32230025:144243 x1,41:18894736,32230025 x1,41:19112968,32230025
k1,42:19112968,32230025:0 ) (1,42:3729360,33016457:15383608,446948,141880
x1,42:5003359,33016457 x1,42:6441853,33016457 k1,41:6776029,33016457:334176
x1,41:7795109,33016457 k1,41:8002621,33016457:207512 x1,41:8293597,33016457
x1,41:8611453,33016457 x1,41:10321924,33016457 k1,41:10529437,33016457:207513
x1,41:11475773,33016457 x1,41:12422747,33016457 k1,41:12630259,33016457:207512
x1,41:13503843,33016457 k1,41:13711355,33016457:207512 x1,41:14948667,33016457
k1,41:15156179,33016457:207512 x1,41:16284698,33016457 x1,41:17740242,33016457
k1,41:17947754,33016457:207512 x1,41:19112968,33016457
k1,41:19112968,33016457:0 ) (1,42:3729360,33802889:15383608,446948,141880
x1,42:4348016,33802889 x1,42:4638992,33802889 x1,42:5331038,33802889
x1,42:5513221,33802889 x1,42:6277356,33802889 k1,41:6525701,33802889:248345
x1,41:7071605,33802889 x1,41:7435971,33802889 x1,41:8964897,33802889
k1,41:9213242,33802889:248345 x1,41:9759154,33802889
k1,41:10007498,33802889:248344 x1,41:10808337,33802889
k1,41:11056682,33802889:248345 x1,41:11566545,33802889 x1,41:12294632,33802889
x1,41:12986678,33802889 x1,41:13678724,33802889 k1,41:13927069,33802889:248345
x1,41:15528748,33802889 x1,41:16530131,33802889 k1,41:16799601,33802889:269470
x1,41:17309464,33802889 k1,41:17557809,33802889:248345 x1,41:17885489,33802889
x1,41:18166641,33802889 x1,41:19112968,33802889 k1,41:19112968,33802889:0 )
(1,42:3729360,34589321:15383608,446948,141880 x1,42:4712382,34589321
k1,41:4890550,34589321:178168 x1,41:5691389,34589321
k1,41:5869557,34589321:178168 x1,41:6925333,34589321
k1,41:7103487,34589321:178154 x1,41:8559686,34589321 x1,41:9397229,34589321
k1,41:9575397,34589321:178168 x1,41:10631172,34589321
k1,41:10809340,34589321:178168 x1,41:11137020,34589321 x1,41:11702588,34589321
x1,41:12867811,34589321 k1,41:13045979,34589321:178168 x1,41:14465474,34589321
x1,41:15302362,34589321 k1,41:15480530,34589321:178168 x1,41:16099186,34589321
x1,41:17191657,34589321 k1,41:17369825,34589321:178168 x1,41:17806936,34589321
k1,41:17985104,34589321:178168 x1,41:18603760,34589321 x1,41:18894736,34589321
x1,41:19112968,34589321 k1,42:19112968,34589321:0 )
(1,42:3729360,35375753:15383608,446948,141880 x1,42:5148855,35375753
k1,41:5315108,35375753:166253 x1,41:5824971,35375753
k1,41:5991224,35375753:166253 x1,41:6901512,35375753 x1,41:8248256,35375753
k1,41:8414509,35375753:166253 x1,41:9215348,35375753
k1,41:9381601,35375753:166253 x1,41:10436712,35375753 x1,41:11492487,35375753
k1,41:11658740,35375753:166253 x1,41:13114939,35375753
k1,41:13281192,35375753:166253 x1,41:15065054,35375753 x1,41:15392734,35375753
x1,41:16653646,35375753 k1,41:16819899,35375753:166253 x1,41:18275443,35375753
x1,41:19112968,35375753 k1,41:19112968,35375753:0 )
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x1,42:4967310,36162185 x1,42:5659356,36162185 x1,42:6351402,36162185
x1,42:6472650,36162185 k1,42:19112968,36162185:12640318
g1,42:19112968,36162185 ) (1,42:3729360,38017577:15383608,541848,14544
k1,42:4878874,38017577:1149514 (1,42:4878874,38017577:0,0,0
g1,42:4878874,38017577 ) x1,42:10209263,38017577 x1,42:11230837,38017577
g1,42:11427445,38017577 x1,42:13742684,38017577 x1,42:14497666,38017577
x1,42:16063444,38017577 g1,42:16260052,38017577 x1,42:17963455,38017577
k1,42:19112968,38017577:1149513 )
(1,42:3729360,38935081:15383608,541848,139977 k1,42:8803160,38935081:5073800
x1,42:12517436,38935081 x1,42:14039169,38935081 g1,42:14039169,38935081
k1,42:19112968,38935081:5073799 )
(1,44:3729360,40160188:15383608,446948,141880
h1,43:3729360,40160188:655361,0,0 x1,43:5403801,40160188
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x1,43:8019615,40160188 k1,43:8305656,40160188:286041 x1,43:8596632,40160188
x1,43:8914488,40160188 x1,43:9533144,40160188 x1,43:9933559,40160188
x1,43:11062078,40160188 x1,43:11717420,40160188 k1,43:12003461,40160188:286041
x1,43:13314163,40160188 x1,43:14224451,40160188 k1,43:14510493,40160188:286042
x1,43:14728725,40160188 (1,43:14728725,40160188:327680,446948,9168
x1,43:15056405,40160188 ) x1,43:15220245,40160188
k1,43:15506286,40160188:286041 (1,43:15506286,40160188:327680,446948,9168
x1,43:15833966,40160188 ) x1,43:15997806,40160188
k1,43:16283847,40160188:286041 (1,43:16283847,40160188:327680,437780,4912
x1,43:16611527,40160188 ) x1,43:16829759,40160188
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k1,43:18603105,40160188:286041 x1,43:19112968,40160188
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x1,44:5149511,40946620 k1,43:5300528,40946620:151017 x1,43:6101367,40946620
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k1,43:12134663,40946620:151016 x1,43:12644526,40946620
k1,43:12795543,40946620:151017 x1,43:13596382,40946620
k1,43:13747399,40946620:151017 x1,43:14839215,40946620 x1,43:16039823,40946620
k1,43:16190840,40946620:151017 x1,43:16736752,40946620
k1,43:16887769,40946620:151017 x1,43:17688608,40946620
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k1,44:19112968,40946620:0 ) (1,44:3729360,41733052:15383608,446948,7531
x1,44:4785135,41733052 g1,43:4948975,41733052 x1,43:6405174,41733052
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g1,44:19112968,41733052 ) (1,46:5040082,43270772:14072886,446948,141880
(1,45:5040082,43270772:0,291271,0 g1,45:5040082,43270772
g1,45:3729360,43270772 g1,45:3401680,43270772
(1,45:3401680,43270772:1310722,291271,0 k1,45:4712402,43270772:1310722
(1,45:4712402,43270772:0,291271,0 k1,45:4384721,43270772:-327681
x1,45:4712402,43270772 ) ) g1,45:5040082,43270772 ) x1,45:5367762,43270772
x1,45:6846906,43270772 x1,45:7720498,43270772 x1,45:8048178,43270772
x1,45:9309090,43270772 k1,45:9711830,43270772:402740 x1,45:10840349,43270772
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x1,45:14992350,43270772 k1,45:15395089,43270772:402739 x1,45:15977697,43270772
x1,45:16815240,43270772 k1,45:17217980,43270772:402740 x1,45:18164316,43270772
k1,45:18567056,43270772:402740 x1,45:19112968,43270772
k1,45:19112968,43270772:0 ) (1,46:5040082,44057204:14072886,446948,141880
x1,46:6386833,44057204 k1,45:6546629,44057204:159796 x1,45:6874309,44057204
x1,45:7962851,44057204 k1,45:8122647,44057204:159796 x1,45:9724983,44057204
x1,45:10107711,44057204 k1,45:10268316,44057204:160605 x1,45:10778179,44057204
k1,45:10937975,44057204:159796 x1,45:12539662,44057204
k1,45:12699458,44057204:159796 x1,45:12990434,44057204
k1,45:13150231,44057204:159797 x1,45:14751918,44057204
k1,45:14911714,44057204:159796 x1,45:16112968,44057204
k1,45:16272764,44057204:159796 x1,45:16818676,44057204
k1,45:16978472,44057204:159796 x1,45:17306152,44057204 x1,45:19112968,44057204
k1,45:19112968,44057204:0 ) (1,46:5040082,44843636:14072886,446948,141880
x1,46:5477193,44843636 k1,45:5674453,44843636:197260 x1,45:6293109,44843636
x1,45:6911765,44843636 x1,45:8440699,44843636 k1,45:8637960,44843636:197261
x1,45:8928936,44843636 x1,45:9246792,44843636 x1,45:10684630,44843636
k1,45:10890245,44843636:205615 x1,45:11436157,44843636
k1,45:11633417,44843636:197260 x1,45:12143280,44843636
k1,45:12340540,44843636:197260 x1,45:13942227,44843636
k1,45:14139488,44843636:197261 x1,45:14649351,44843636 x1,45:14977031,44843636
x1,45:15159214,44843636 x1,45:15980373,44843636 x1,45:16781212,44843636
k1,45:16978472,44843636:197260 x1,45:17306152,44843636 x1,45:19112968,44843636
k1,45:19112968,44843636:0 ) (1,46:5040082,45630068:14072886,446948,92400
x1,46:5840266,45630068 g1,45:6004106,45630068 x1,45:6295082,45630068
x1,45:6612938,45630068 x1,45:8069119,45630068 k1,46:19112968,45630068:11043849
g1,46:19112968,45630068 ) (1,36:3729360,47050772:15383608,435814,113504
(1,36:3729360,47050772:943720,367985,0 k1,36:4345400,47050772:616040
(1,36:4345400,47050772:327680,367985,0
$1,36:4345400,47050772(1,36:4345400,46865742:283988,182955,0x1,36:4596620,46865742)$1,36:4629388,47050772
k1,36:4673080,47050772:43692 ) ) (1,36:4673080,47050772:0,435814,0
r1,47:4673080,47050772:0,435814,0 ) x1,36:5168005,47050772
x1,36:5824939,47050772 k1,36:6103556,47050772:278617 x1,36:6831787,47050772
x1,36:8171333,47050772 k1,36:8449951,47050772:278618 x1,36:8974239,47050772
k1,36:9252856,47050772:278617 x1,36:9893528,47050772
k1,36:10172145,47050772:278617 x1,36:10842180,47050772
k1,36:11120649,47050772:278469 x1,36:11731964,47050772 x1,36:12926274,47050772
k1,36:13204891,47050772:278617 x1,36:13874920,47050772 x1,36:14835407,47050772
k1,36:15114025,47050772:278618 x1,36:15405526,47050772 x1,36:16184099,47050772
x1,36:16679024,47050772 x1,36:17494806,47050772 k1,36:17773423,47050772:278617
x1,36:18996578,47050772 x1,36:19112968,47050772 k1,36:19112968,47050772:0 )
(1,36:3729360,47673364:15383608,357558,113504 x1,36:5695440,47673364
k1,36:5861950,47673364:166510 x1,36:6910519,47673364 x1,36:7026909,47673364
x1,36:8992989,47673364 k1,36:9159498,47673364:166509 x1,36:9916567,47673364
k1,36:10075990,47673364:159423 x1,36:11124559,47673364 x1,36:11240949,47673364
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x1,36:16852689,47673364 k1,36:17012112,47673364:159423 x1,36:17448842,47673364
k1,36:17608264,47673364:159422 x1,36:18190747,47673364 x1,36:18598114,47673364
x1,36:19112968,47673364 k1,36:19112968,47673364:0 )
(1,36:3729360,48295956:15383608,357558,186779 x1,36:4777405,48295956
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x1,36:13471621,48295956 x1,36:13995909,48295956 x1,36:14228690,48295956
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x1,36:16311142,48295956 x1,36:17388544,48295956
r1,47:17388544,48295956:0,186779,186779 k1,36:19112968,48295956:1724424
g1,36:19112968,48295956 ) ] g1,97:19112968,48295956 )
k1,97:19579138,48295956:466170 r1,97:19579138,48295956:0,33750705,0
k1,97:20045308,48295956:466170 (1,97:20045308,48295956:15383608,33750705,0
[1,97:20045308,48295956:15383608,33750705,0
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(1,46:21028350,15200611:0,291271,0 k1,46:20700669,15200611:-327681
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(1,48:21028350,20890065:0,291271,0 k1,48:20700669,20890065:-327681
x1,48:21028350,20890065 ) ) g1,48:21356030,20890065 ) x1,48:23649772,20890065
k1,48:23852359,20890065:202587 x1,48:24834726,20890065
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k1,48:34919053,20890065:202587 x1,48:35428916,20890065
k1,48:35428916,20890065:0 ) (1,49:21356030,21676497:14072886,446948,141880
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k1,48:31392576,21676497:218314 x1,48:33267543,21676497
k1,48:33499475,21676497:231932 x1,48:34919626,21676497
k1,48:35137940,21676497:218314 x1,48:35428916,21676497
k1,48:35428916,21676497:0 ) (1,49:21356030,22462929:14072886,446948,141880
x1,49:22739485,22462929 x1,49:24414546,22462929 x1,48:25417888,22462929
x1,48:25539136,22462929 k1,48:25792203,22462929:253067 x1,48:27648809,22462929
k1,48:27884031,22462929:235222 x1,48:28830367,22462929
k1,48:29065589,22462929:235222 x1,48:29393269,22462929 x1,48:30481811,22462929
k1,48:30717034,22462929:235223 x1,48:32264978,22462929
k1,48:32500200,22462929:235222 x1,48:32791176,22462929
k1,48:33026398,22462929:235222 x1,48:35428916,22462929
k1,48:35428916,22462929:0 ) (1,49:21356030,23249361:14072886,446948,141880
x1,49:23575709,23249361 x1,48:23713340,23249361 k1,48:23934774,23249361:221434
x1,48:24225750,23249361 k1,48:24435665,23249361:209915 x1,48:26037344,23249361
x1,48:27082632,23249361 k1,48:27292547,23249361:209915 x1,48:28421066,23249361
x1,48:28923721,23249361 k1,48:29133636,23249361:209915 x1,48:30079972,23249361
k1,48:30289888,23249361:209916 x1,48:30908544,23249361
k1,48:31118459,23249361:209915 x1,48:32138187,23249361 x1,48:32720139,23249361
x1,48:33726761,23249361 k1,48:33936676,23249361:209915 x1,48:35428916,23249361
k1,48:35428916,23249361:0 ) (1,49:21356030,24035793:14072886,446948,7531
x1,49:24050187,24035793 x1,48:24177987,24035793
k1,49:35428916,24035793:11250929 g1,49:35428916,24035793 )
(1,51:20045308,25268798:15383608,446948,141880 x1,51:21064388,25268798
k1,50:21241641,25268798:177253 x1,50:22151911,25268798
k1,50:22329163,25268798:177252 x1,50:23420978,25268798 x1,50:24877177,25268798
k1,50:25054430,25268798:177253 x1,50:25491541,25268798
k1,50:25668793,25268798:177252 x1,50:26652453,25268798
k1,50:26829706,25268798:177253 x1,50:27448362,25268798 x1,50:28285905,25268798
k1,50:28463158,25268798:177253 x1,50:28790838,25268798 x1,50:29081814,25268798
x1,50:29393110,25268798 x1,50:29856445,25268798 x1,50:31294301,25268798
k1,50:31474906,25268798:180605 x1,50:32421242,25268798
k1,50:32598495,25268798:177253 x1,50:33618212,25268798
k1,50:33795464,25268798:177252 x1,50:34633007,25268798
k1,50:34810260,25268798:177253 x1,50:35428916,25268798
k1,50:35428916,25268798:0 ) (1,51:20045308,26055230:15383608,446948,141880
x1,51:20810099,26055230 x1,51:22557267,26055230 k1,50:22779956,26055230:222689
x1,50:23325868,26055230 x1,50:24563171,26055230 k1,50:24785860,26055230:222689
x1,50:25295723,26055230 k1,50:25518412,26055230:222689 x1,50:26465386,26055230
k1,50:26688075,26055230:222689 x1,50:27306731,26055230 x1,50:28143619,26055230
x1,50:28271419,26055230 k1,50:28651126,26055230:379707 x1,50:29051541,26055230
k1,50:29274230,26055230:222689 x1,50:30183862,26055230
k1,50:30406551,26055230:222689 x1,50:31061911,26055230 x1,50:31862750,26055230
x1,50:33210138,26055230 k1,50:33432827,26055230:222689 x1,50:34051483,26055230
k1,50:34274172,26055230:222689 x1,50:34601852,26055230 x1,50:35428916,26055230
k1,50:35428916,26055230:0 ) (1,51:20045308,26841662:15383608,446948,141880
x1,51:20918892,26841662 x1,51:21810827,26841662 k1,50:22116148,26841662:305321
x1,50:23462891,26841662 k1,50:23739917,26841662:277026 x1,50:24831732,26841662
x1,50:25923548,26841662 x1,50:27489178,26841662 k1,50:27766203,26841662:277025
x1,50:28312115,26841662 k1,50:28589141,26841662:277026 x1,50:28916821,26841662
x1,50:30395965,26841662 x1,50:31269557,26841662 x1,50:31597237,26841662
x1,50:32858149,26841662 k1,50:33135174,26841662:277025 x1,50:34591373,26841662
x1,50:35428916,26841662 k1,50:35428916,26841662:0 )
(1,51:20045308,27628094:15383608,446948,141880 x1,51:21282603,27628094
k1,50:21532275,27628094:249672 x1,50:22078187,27628094 x1,50:22405867,27628094
x1,50:23524562,27628094 k1,50:23774233,27628094:249671 x1,50:24392889,27628094
x1,50:24902752,27628094 x1,50:26395637,27628094 k1,50:26645309,27628094:249672
x1,50:27155172,27628094 x1,50:27701084,27628094 x1,50:28720155,27628094
x1,50:29739881,27628094 k1,50:29989553,27628094:249672 x1,50:30644913,27628094
k1,50:30894584,27628094:249671 x1,50:31695423,27628094
k1,50:31945095,27628094:249672 x1,50:32563751,27628094 x1,50:33582822,27628094
k1,50:33832493,27628094:249671 x1,50:34378405,27628094
k1,50:34628077,27628094:249672 x1,50:35428916,27628094
k1,50:35428916,27628094:0 ) (1,51:20045308,28414526:15383608,446948,141880
x1,51:20336284,28414526 x1,51:21338316,28414526 x1,51:21702682,28414526
x1,51:22013978,28414526 x1,51:22295130,28414526 k1,50:22453541,28414526:158411
x1,50:25292531,28414526 k1,50:25450942,28414526:158411 x1,50:25996854,28414526
k1,50:26155266,28414526:158412 x1,50:26956105,28414526
k1,50:27114516,28414526:158411 x1,50:28570715,28414526
k1,50:28729126,28414526:158411 x1,50:30330813,28414526
k1,50:30489224,28414526:158411 x1,50:31398856,28414526 x1,50:32236399,28414526
k1,50:32394811,28414526:158412 x1,50:33013467,28414526
k1,50:33171878,28414526:158411 x1,50:33827238,28414526 x1,50:35428916,28414526
k1,50:35428916,28414526:0 ) (1,51:20045308,29200958:15383608,446948,141880
x1,51:20555171,29200958 g1,50:20719011,29200958 x1,50:21337667,29200958
x1,50:21655523,29200958 g1,50:21819363,29200958 x1,50:23166114,29200958
g1,50:23329954,29200958 x1,50:23803122,29200958 x1,50:24415225,29200958
x1,50:24536473,29200958 k1,51:35428916,29200958:10892443
g1,51:35428916,29200958 ) (1,52:20045308,30699541:15383608,541848,14544
k1,52:20833478,30699541:788170 (1,52:20833478,30699541:0,0,0
g1,52:20833478,30699541 ) x1,52:25334185,30699541 g1,52:25530793,30699541
x1,52:27715483,30699541 g1,52:27912091,30699541 x1,52:34640747,30699541
g1,52:34640747,30699541 k1,52:35428916,30699541:788169 )
(1,65:20045308,31780884:15383608,446948,141880
h1,63:20045308,31780884:655361,0,0 x1,63:21173837,31780884
k1,63:21453041,31780884:279204 x1,63:22144432,31780884 x1,63:22751296,31780884
k1,63:23030616,31780884:279320 x1,63:23358296,31780884 x1,63:23632888,31780884
x1,63:24033303,31780884 x1,63:24834142,31780884 k1,63:25113463,31780884:279321
x1,63:25659375,31780884 k1,63:25938695,31780884:279320 x1,63:26557351,31780884
x1,63:27176007,31780884 x1,63:28195078,31780884 k1,63:28474399,31780884:279321
x1,63:29602918,31780884 k1,63:29882238,31780884:279320 x1,63:31047444,31780884
k1,63:31326764,31780884:279320 x1,63:32418579,31780884
k1,63:32697900,31780884:279321 x1,63:33316556,31780884 x1,63:33631132,31780884
x1,63:33912284,31780884 k1,63:34191604,31780884:279320 x1,63:35428916,31780884
k1,63:35428916,31780884:0 ) (1,65:20045308,32567316:15383608,446948,141880
x1,65:20846147,32567316 x1,63:21622090,32567316 k1,63:21886153,32567316:264063
x1,63:22504809,32567316 x1,63:22822665,32567316 x1,63:23103817,32567316
x1,63:24159592,32567316 k1,63:24403611,32567316:244019 x1,63:24694587,32567316
k1,63:24938606,32567316:244019 x1,63:26212613,32567316
k1,63:26456632,32567316:244019 x1,63:27912176,32567316
k1,63:28156195,32567316:244019 x1,63:28702107,32567316
k1,63:28946126,32567316:244019 x1,63:30365621,32567316
k1,63:30609640,32567316:244019 x1,63:31555976,32567316
k1,63:31799995,32567316:244019 x1,63:32600834,32567316 x1,63:34020329,32567316
k1,63:34264348,32567316:244019 x1,63:35210684,32567316 x1,63:35428916,32567316
k1,65:35428916,32567316:0 ) (1,65:20045308,33353748:15383608,446948,141880
x1,65:20555171,33353748 x1,65:21993665,33353748 k1,63:22192125,33353748:198460
x1,63:23211205,33353748 k1,64:23360944,33353748:149739 x1,64:24307280,33353748
x1,64:24889878,33353748 x1,64:25217558,33353748 x1,64:25532134,33353748
x1,64:26551860,33353748 k1,64:26701598,33353748:149738 x1,64:27211461,33353748
k1,64:27361200,33353748:149739 x1,64:28162039,33353748
k1,64:28311778,33353748:149739 x1,64:28639458,33353748 x1,64:30118602,33353748
x1,64:30992194,33353748 x1,64:31319874,33353748 x1,64:32580786,33353748
k1,64:32730524,33353748:149738 x1,64:33859043,33353748
k1,64:34008782,33353748:149739 x1,64:35428916,33353748
k1,64:35428916,33353748:0 ) (1,65:20045308,34140180:15383608,446948,141880
x1,65:21319315,34140180 k1,64:21552521,34140180:233206 x1,64:22681031,34140180
k1,64:22914236,34140180:233205 x1,64:23860572,34140180
k1,64:24093778,34140180:233206 x1,64:25732160,34140180
k1,64:25965365,34140180:233205 x1,64:27384860,34140180 x1,64:28367882,34140180
k1,64:28601088,34140180:233206 x1,64:29074247,34140180 x1,64:29365223,34140180
x1,64:30930853,34140180 k1,64:31164058,34140180:233205 x1,64:31782714,34140180
x1,64:32097290,34140180 x1,64:32378442,34140180 k1,64:32611648,34140180:233206
x1,64:33848960,34140180 k1,64:34082165,34140180:233205 x1,64:35210684,34140180
x1,64:35428916,34140180 k1,65:35428916,34140180:0 )
(1,65:20045308,34926612:15383608,446948,141880 x1,65:21501507,34926612
k1,64:21691401,34926612:189894 x1,64:22965408,34926612
k1,64:23155302,34926612:189894 x1,64:23956141,34926612
k1,64:24146035,34926612:189894 x1,64:25201146,34926612 x1,64:26074720,34926612
k1,64:26264614,34926612:189894 x1,64:27429829,34926612
k1,64:27619722,34926612:189893 x1,64:28456610,34926612 x1,64:28966473,34926612
x1,64:30022240,34926612 k1,64:30212134,34926612:189894 x1,64:30758046,34926612
k1,64:30947940,34926612:189894 x1,64:31748779,34926612
k1,64:31938673,34926612:189894 x1,64:33066537,34926612 x1,64:33673401,34926612
k1,64:33863295,34926612:189894 x1,64:34628086,34926612 x1,64:35210684,34926612
x1,64:35428916,34926612 k1,65:35428916,34926612:0 )
(1,65:20045308,35713044:15383608,446948,141880 x1,65:20372988,35713044
x1,65:20687564,35713044 x1,65:21962218,35713044 k1,64:22188108,35713044:225890
x1,64:23024996,35713044 x1,64:24626674,35713044 k1,64:24852564,35713044:225890
x1,64:25507924,35713044 k1,64:25733814,35713044:225890 x1,64:26534653,35713044
k1,64:26760542,35713044:225889 x1,64:27815653,35713044 x1,64:28871428,35713044
k1,64:29097318,35713044:225890 x1,64:30553517,35713044
k1,64:30779407,35713044:225890 x1,64:31726381,35713044 x1,64:32054061,35713044
x1,64:32527220,35713044 x1,64:33546946,35713044 k1,64:33772836,35713044:225890
x1,64:35428916,35713044 k1,64:35428916,35713044:0 )
(1,65:20045308,36499476:15383608,446948,141880 x1,65:21064388,36499476
k1,64:21284076,36499476:219688 x1,64:22339851,36499476 x1,64:23613858,36499476
k1,64:23833546,36499476:219688 x1,64:24124522,36499476 x1,64:24439098,36499476
x1,64:24720250,36499476 x1,64:25848114,36499476 k1,64:26067802,36499476:219688
x1,64:27196321,36499476 k1,64:27416010,36499476:219689 x1,64:28836144,36499476
k1,64:29055832,36499476:219688 x1,64:30329184,36499476
k1,64:30548872,36499476:219688 x1,64:31677391,36499476
k1,64:31897079,36499476:219688 x1,64:32661214,36499476 x1,64:33279870,36499476
x1,64:34553868,36499476 k1,64:34773556,36499476:219688 x1,64:35428916,36499476
k1,64:35428916,36499476:0 ) (1,65:20045308,37285908:15383608,446948,141880
x1,65:20846147,37285908 x1,65:21683690,37285908 k1,64:21930657,37285908:246967
x1,64:22221633,37285908 k1,64:22468600,37285908:246967 x1,64:23778647,37285908
x1,64:26181173,37285908 k1,64:26428140,37285908:246967 x1,64:26719116,37285908
x1,64:27033692,37285908 x1,64:27314844,37285908 x1,64:28442708,37285908
k1,64:28689676,37285908:246968 x1,64:29235588,37285908
k1,64:29482555,37285908:246967 x1,64:30283394,37285908
k1,64:30530361,37285908:246967 x1,64:31695576,37285908
k1,64:31942543,37285908:246967 x1,64:32779431,37285908 x1,64:33289294,37285908
x1,64:34345061,37285908 k1,64:34592028,37285908:246967 x1,64:35428916,37285908
k1,64:35428916,37285908:0 ) (1,65:20045308,38072340:15383608,446948,141880
x1,65:21464803,38072340 x1,65:22447825,38072340 g1,64:22611665,38072340
x1,64:23084824,38072340 x1,64:23375800,38072340 x1,64:24850342,38072340
g1,64:25014182,38072340 x1,64:25960518,38072340 g1,64:26124358,38072340
x1,64:27252877,38072340 x1,64:28708421,38072340 g1,64:28872261,38072340
x1,64:30037475,38072340 g1,64:30201315,38072340 x1,64:31293130,38072340
x1,64:32384946,38072340 x1,64:34114416,38072340
k1,65:35428916,38072340:1314500 g1,65:35428916,38072340 )
(1,77:20045308,38858772:15383608,446948,141880
h1,75:20045308,38858772:655361,0,0 x1,75:22156859,38858772
k1,75:22456624,38858772:299765 x1,75:23767344,38858772
k1,75:24067109,38858772:299765 x1,75:25485949,38858772 x1,75:26796651,38858772
k1,75:27096415,38858772:299764 x1,75:27314647,38858772
(1,75:27314647,38858772:327680,446948,9168 x1,75:27642327,38858772 )
x1,75:27806167,38858772 k1,75:28105932,38858772:299765
(1,75:28105932,38858772:327680,446948,14416 x1,75:28433612,38858772 )
x1,75:28597452,38858772 k1,75:28897217,38858772:299765
(1,75:28897217,38858772:655360,446948,9168 x1,75:29552577,38858772 )
x1,75:29716417,38858772 k1,75:30016182,38858772:299765
(1,75:30016182,38858772:655360,446948,0 x1,75:30671542,38858772 )
x1,75:30889774,38858772 k1,75:31189538,38858772:299764 x1,75:32063122,38858772
k1,75:32362887,38858772:299765 x1,75:33600199,38858772
k1,75:33899964,38858772:299765 x1,75:35428916,38858772
k1,75:35428916,38858772:0 ) (1,77:20045308,39645204:15383608,446948,141880
x1,77:20663964,39645204 x1,77:22228947,39645204 k1,75:22458302,39645204:229355
x1,75:23259141,39645204 k1,75:23488495,39645204:229354 x1,75:24580310,39645204
x1,75:25672126,39645204 x1,75:27237756,39645204 k1,75:27467111,39645204:229355
x1,75:28413447,39645204 k1,75:28642802,39645204:229355 x1,75:29443641,39645204
k1,75:29672996,39645204:229355 x1,75:30801515,39645204 x1,75:32257714,39645204
k1,75:32487068,39645204:229354 x1,75:33615587,39645204
k1,75:33844942,39645204:229355 x1,75:35428916,39645204
k1,75:35428916,39645204:0 ) (1,77:20045308,40431636:15383608,446948,141880
x1,77:20627916,40431636 x1,77:21465459,40431636 k1,76:21915010,40431636:449551
x1,76:22388169,40431636 x1,76:23407240,40431636 k1,76:23856792,40431636:449552
x1,76:24402704,40431636 k1,76:24852255,40431636:449551 x1,76:25653094,40431636
k1,76:26102646,40431636:449552 x1,76:26430326,40431636 x1,76:28092308,40431636
k1,76:28541859,40431636:449551 x1,76:30034099,40431636
k1,76:30483650,40431636:449551 x1,76:32892077,40431636
k1,76:33341629,40431636:449552 x1,76:33559861,40431636
(1,76:33559861,40431636:655360,446948,0 x1,76:34215221,40431636 )
x1,76:34433453,40431636 k1,76:34883004,40431636:449551 x1,76:35428916,40431636
k1,77:35428916,40431636:0 ) (1,77:20045308,41218068:15383608,446948,141880
x1,77:22109674,41218068 k1,76:22274547,41218068:164873 x1,76:22492779,41218068
(1,76:22492779,41218068:655360,446948,9168 x1,76:23148139,41218068 )
x1,76:23530211,41218068 k1,76:23736471,41218068:206260 x1,76:24901686,41218068
k1,76:25066560,41218068:164874 x1,76:26485400,41218068 x1,76:27959942,41218068
k1,76:28125073,41218068:165131 x1,76:29726760,41218068
k1,76:29891634,41218068:164874 x1,76:30328745,41218068
k1,76:30493618,41218068:164873 x1,76:32022570,41218068
k1,76:32187443,41218068:164873 x1,76:32951587,41218068
k1,76:33116461,41218068:164874 x1,76:33917300,41218068
k1,76:34082173,41218068:164873 x1,76:35428916,41218068
k1,76:35428916,41218068:0 ) (1,77:20045308,42004500:15383608,446948,141880
x1,77:20773395,42004500 k1,76:20964828,42004500:191433 x1,76:21510740,42004500
k1,76:21702172,42004500:191432 x1,76:22649146,42004500 x1,76:22976826,42004500
x1,76:23449985,42004500 x1,76:24469711,42004500 k1,76:24661144,42004500:191433
x1,76:25279800,42004500 x1,76:26298871,42004500 x1,76:26589847,42004500
x1,76:27536166,42004500 k1,76:27727599,42004500:191433 x1,76:28346255,42004500
x1,76:29656957,42004500 k1,76:29848389,42004500:191432 x1,76:31486109,42004500
k1,76:31677542,42004500:191433 x1,76:31968518,42004500
k1,76:32159950,42004500:191432 x1,76:33433957,42004500
k1,76:33625390,42004500:191433 x1,76:33953070,42004500 x1,76:34227662,42004500
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k1,78:35428916,48295956:0 ) ] g1,97:35428916,48295956 ) ) ]
h1,97:3729360,50262036:31699556,0,0 ] ] !58819 1 !10 2
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k1,78:17984458,24897538:183576 x1,78:19112968,24897538
k1,78:19112968,24897538:0 ) ] g1,103:19112968,25521564 )
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k1,103:20045308,25521564:466170 (1,103:20045308,25521564:15383608,6784410,0
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x1,101:35428916,25521564 k1,101:35428916,25521564:0 ) ]
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x1,71:8170703,48154077 k1,71:35428916,48154077:27258213
g1,71:35428916,48154077 ) ] ] h1,103:3729360,50262038:31699556,0,0 ] ] !18074
2 !10 3 [1,135:1,50262038:35428915,50262037,0
[1,135:3729360,50262038:31699556,49329697,0
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k1,105:12762261,32851029:161514
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k1,105:14792649,32851029:161514
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x1,105:15689296,32851029 k1,105:15851275,32851029:161979
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x1,105:18200705,32851029 x1,105:18991720,32851029 x1,105:19112968,32851029
k1,105:19112968,32851029:0 ) ] g1,135:19112968,32851029 )
k1,135:19579138,32851029:466170 r1,135:19579138,32851029:0,15765882,0
k1,135:20045308,32851029:466170 (1,135:20045308,32851029:15383608,15765882,0
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(1,90:4095355,46098218:40258437,31377900,0 ) k1,90:4095355,46098218:-40258437
) ) ) )
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x1,95:5731471,48288426 x1,95:5859271,48288426 k1,95:35428916,48288426:29569645
g1,95:35428916,48288426 ) ] ] h1,135:3729360,50262038:31699556,0,0 ] ] !34569
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g1,178:35428916,48295958 ) ) ] h1,178:3729360,50262038:31699556,0,0 ] ] !13498
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(1,203:20159061,37762657:13090570,9867264,0 )
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k1,204:35428916,41563745:7710691 g1,204:35428916,41563745 ) ]
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x1,245:35428916,48295956 k1,245:35428916,48295956:0 ) ]
g1,246:35428916,48295956 ) ) ] h1,246:3729360,50262036:31699556,0,0 ] ] !44828
6 !10 7 [1,311:1,50262037:35428915,50262036,0
[1,311:3729360,50262037:31699556,49329696,0
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(1,254:5137229,26138620:0,0,0 g1,254:5137229,26138620 )
x1,254:7539762,26138620 x1,254:10977221,26138620 g1,254:11173829,26138620
x1,254:13402555,26138620 g1,254:13599163,26138620 x1,254:14734765,26138620
x1,254:15811385,26138620 x1,254:17705099,26138620
k1,254:19112968,26138620:1407869 )
(1,254:3729360,27056124:15383608,541848,14544 k1,254:5356651,27056124:1627291
x1,254:8459092,27056124 x1,254:9870723,27056124 g1,254:10067331,27056124
x1,254:11770734,27056124 g1,254:11967342,27056124 x1,254:15113040,27056124
x1,254:17485678,27056124 g1,254:17485678,27056124
k1,254:19112968,27056124:1627290 )
(1,280:3729360,28296037:15383608,538305,162393 (1,280:3729360,28296037:0,0,0
g1,280:3729360,28296037 ) x1,280:6350503,28296037 g1,280:6547111,28296037
x1,280:9474158,28296037 k1,280:19112968,28296037:9638810
g1,280:19112968,28296037 ) (1,282:3729360,29343942:15383608,446948,141880
h1,281:3729360,29343942:655361,0,0 x1,281:5403801,29343942
k1,281:5625368,29343942:221567 x1,281:6244024,29343942 x1,281:6535000,29343942
x1,281:7700205,29343942 k1,281:7921772,29343942:221567 x1,281:8722611,29343942
x1,281:9596203,29343942 k1,281:9817771,29343942:221568
x1,281:10582562,29343942 x1,281:12329730,29343942
k1,281:12551297,29343942:221567 x1,281:13061160,29343942
k1,281:13282727,29343942:221567 x1,281:13500959,29343942
(1,281:13500959,29343942:655360,446948,14416 x1,281:14156319,29343942 )
x1,281:14374551,29343942 k1,281:14596118,29343942:221567
x1,281:15069286,29343942 x1,281:15608637,29343942
k1,281:15830205,29343942:221568 x1,281:17031469,29343942
k1,281:17253036,29343942:221567 x1,281:17762899,29343942
k1,281:17984466,29343942:221567 x1,281:18530370,29343942
x1,281:18894736,29343942 x1,281:19112968,29343942 k1,282:19112968,29343942:0 )
(1,282:3729360,30130374:15383608,446948,141880 x1,282:5003358,30130374
k1,281:5133352,30130374:129994 x1,281:5934191,30130374
k1,281:6064186,30130374:129995 x1,281:7519730,30130374 x1,281:8357255,30130374
k1,281:8487249,30130374:129994 x1,281:8997112,30130374 x1,281:9725199,30130374
x1,281:10417245,30130374 x1,281:11109291,30130374
k1,281:11239285,30130374:129994 x1,281:12840964,30130374
x1,281:13678507,30130374 k1,281:13808502,30130374:129995
x1,281:14318365,30130374 k1,281:14448359,30130374:129994
x1,281:15249198,30130374 k1,281:15379193,30130374:129995
x1,281:16725936,30130374 x1,281:17563479,30130374
k1,281:17693473,30130374:129994 x1,281:19112968,30130374
k1,281:19112968,30130374:0 ) (1,282:3729360,30916806:15383608,446948,141880
x1,282:4348016,30916806 x1,282:5549288,30916806 x1,282:5731471,30916806
x1,282:6532310,30916806 x1,282:7715876,30916806 k1,281:8131323,30916806:415447
x1,281:8677235,30916806 k1,281:8911837,30916806:234602 x1,281:9858811,30916806
k1,281:10093413,30916806:234602 x1,281:11895635,30916806
k1,281:12147927,30916806:252292 x1,281:12948766,30916806
k1,281:13183369,30916806:234603 x1,281:14238480,30916806
x1,281:15294255,30916806 k1,281:15528857,30916806:234602
x1,281:16985056,30916806 k1,281:17219658,30916806:234602
x1,281:18165994,30916806 x1,281:19112968,30916806 k1,281:19112968,30916806:0 )
] g1,311:19112968,30916806 ) k1,311:19579138,30916806:466170
r1,311:19579138,30916806:0,27559633,0 k1,311:20045308,30916806:466170
(1,311:20045308,30916806:15383608,27559633,0
[1,311:20045308,30916806:15383608,27559633,0
(1,282:20045308,4012533:15383608,446948,141880 x1,282:20663964,4012533
x1,282:21391404,4012533 x1,282:22957052,4012533 x1,282:23794595,4012533
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x1,281:26834767,4012533 x1,281:27526813,4012533 x1,281:28218859,4012533
g1,281:28382699,4012533 x1,281:29984378,4012533 x1,281:30821921,4012533
g1,281:30985761,4012533 x1,281:31422872,4012533 g1,281:31586712,4012533
x1,281:32096575,4012533 x1,281:32407871,4012533 x1,281:33016703,4012533
g1,281:33180543,4012533 x1,281:33835903,4012533
k1,282:35428916,4012533:1593013 g1,282:35428916,4012533 )
(1,283:20045308,4798965:15383608,0,0 h1,283:20045308,4798965:655361,0,0
k1,283:35428916,4798965:14728247 g1,283:35428916,4798965 )
(1,287:24239379,5730200:11189537,865699,709964
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(1,287:24578210,5828503:941864,318577,89201 x1,287:25487306,5828503 )
g1,287:25702114,5730200 x1,287:26211840,5730200 g1,287:26393880,5730200
(1,287:26393880,5730200:4658920,865699,709964
(1,287:26393880,5730200:4658920,865699,709964
h1,287:26393880,5730200:78643,0,0
[1,287:26472523,5730200:4501634,865699,709964
(1,287:26472523,5286844:4501634,422343,127431 k1,287:27864690,5286844:1392167
x1,287:28531671,5286844 x1,287:28891815,5286844
(1,287:28891815,5385147:280122,197518,0 x1,287:29139169,5385147 )
x1,287:29581990,5286844 k1,287:30974157,5286844:1392167 )
(1,287:26472523,6252659:4501634,594554,187504
(1,287:26472523,6252659:873816,594554,87018
(1,287:26472523,5710531:546135,26213,629146 x1,287:27018658,5710531 )
[1,287:27018658,6252659:327681,594554,0
(1,287:27018658,6252659:327681,422343,0 x1,287:27346339,6252659 ) ] )
x1,287:27902940,6252659 x1,287:28763585,6252659
(1,287:28763585,6350962:281419,197518,0 x1,287:29012236,6350962 )
g1,287:29154228,6252659 (1,287:29154228,6252659:1195429,491520,187504
(1,287:29154228,6252659:254863,491520,163840 x1,287:29409091,6252659 )
x1,287:29779779,6252659 (1,287:29779779,6350962:315015,197518,89201
x1,287:30045566,6350962 ) (1,287:30094794,6252659:254863,491520,163840
x1,287:30349657,6252659 ) ) g1,287:30458881,6252659 x1,287:30927281,6252659 )
] h1,287:30974157,5730200:78643,0,0 ) ) x1,287:31234845,5730200 )
(1,287:34664772,5730200:764144,446948,116320
(1,287:34664772,5730200:764144,446948,116320 x1,287:35428916,5730200 ) ) )
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k1,288:23880421,7297468:170649 x1,288:24681260,7297468
k1,288:24851908,7297468:170648 x1,288:27982539,7297468
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k1,288:29617898,7297468:172351
$1,288:29617898,7297468x1,288:30174499,7297468$1,288:30221375,7297468
k1,288:30392023,7297468:170648 x1,288:30829134,7297468
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k1,288:31971270,7297468:170648 x1,288:35428916,7297468
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x1,290:20263540,8083900 x1,288:21567034,8083900 k1,288:21733987,8083900:166953
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k1,288:32225674,8083900:167730 x1,288:33172010,8083900
k1,288:33338963,8083900:166953
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k1,288:34024013,8083900:166952 x1,288:34461124,8083900
k1,288:34628077,8083900:166953 x1,288:35428916,8083900
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k1,288:28280063,8870332:212895 x1,288:29445277,8870332
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k1,288:33318030,8870332:212895 x1,288:35428916,8870332
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x1,290:22776143,9656764 k1,288:22932442,9656764:156299 x1,288:26116801,9656764
k1,288:26317448,9656764:200647
$1,288:26317448,9656764x1,288:26727501,9656764$1,288:26727501,9656764
k1,288:26883800,9656764:156299 x1,288:27320911,9656764
k1,288:27477210,9656764:156299 x1,288:28278049,9656764
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k1,289:33318030,9656764:156299 x1,289:35428916,9656764
k1,289:35428916,9656764:0 ) (1,290:20045308,10443196:15383608,446948,7531
x1,290:22776143,10443196 g1,289:22939983,10443196 x1,289:23449846,10443196
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g1,289:26617172,10443196 x1,289:27418011,10443196 g1,289:27581851,10443196
x1,289:27909531,10443196 x1,289:29352626,10443196
k1,290:35428916,10443196:6076290 g1,290:35428916,10443196 )
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(1,293:26023670,11802861:3426885,725524,449545 x1,293:26433723,11802861
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(1,293:27307529,11802861:1960981,725524,449545
(1,293:27307529,11802861:1960981,725524,449545
h1,293:27307529,11802861:78643,0,0
[1,293:27386172,11802861:1803695,725524,449545
(1,293:27386172,11359505:1803695,282168,187504 x1,293:27794103,11359505
(1,293:27794103,11457808:719314,197518,89201 x1,293:28464189,11457808 )
x1,293:28887898,11359505 (1,293:28887898,11457808:301969,197518,0
x1,293:29138288,11457808 ) ) (1,293:27386172,12252405:1803695,282168,0
k1,293:28146214,12252405:760042 x1,293:28429825,12252405
k1,293:29189867,12252405:760042 ) ] h1,293:29189867,11802861:78643,0,0 ) )
x1,293:29450555,11802861 ) (1,293:34664772,11802861:764144,446948,116320
(1,293:34664772,11802861:764144,446948,116320 x1,293:35428916,11802861 ) ) )
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k1,294:21828615,13338124:182275
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x1,294:27862928,13338124 k1,294:28049811,13338124:186883
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x1,295:21100419,14124556 x1,295:22156194,14124556
k1,294:22402413,14124556:246219 x1,294:23239956,14124556
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h1,296:20045308,15697420:655361,0,0 x1,296:22630036,15697420
k1,296:22907324,15697420:277288 x1,296:23562684,15697420
k1,296:23839971,15697420:277287 x1,296:24495331,15697420
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$1,297:30229964,15697420x1,297:30786565,15697420$1,297:30833441,15697420
k1,297:31110728,15697420:277287 x1,297:31547839,15697420
k1,297:31825127,15697420:277288 x1,297:32443783,15697420
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$1,297:20836268,16483852x1,297:21246321,16483852$1,297:21246321,16483852
x1,297:21410161,16483852 k1,297:21675521,16483852:265360
x1,297:22621857,16483852 k1,297:22866914,16483852:245057
$1,297:22866914,16483852x1,297:23385012,16483852$1,297:23385012,16483852
k1,297:23630068,16483852:245056 x1,297:24067179,16483852
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k1,297:32888408,16483852:245057 x1,297:33834744,16483852
k1,297:34079800,16483852:245056
$1,297:34079800,16483852x1,297:34489853,16483852$1,297:34489853,16483852
x1,297:34653693,16483852 k1,297:34919053,16483852:265360
x1,297:35428916,16483852 k1,297:35428916,16483852:0 )
(1,298:20045308,17270284:15383608,446948,141880 x1,298:21646995,17270284
k1,297:21805249,17270284:158254 x1,297:22606088,17270284
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h1,311:3729360,50262037:31699556,0,0 ] ] !57182 7 !10 8
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x1,322:35428916,48295958 k1,323:35428916,48295958:0 ) ]
g1,323:35428916,48295958 ) ) ] h1,323:3729360,50262038:31699556,0,0 ] ] !35383
8 !10 9 [1,386:1,50262036:35428915,50262035,0
[1,386:3729360,50262036:31699556,49329695,0
[1,386:3729360,1718773:31699556,786432,0 h1,386:3729360,1718773:31699556,0,0 ]
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x1,341:10137934,28812107 x1,341:12074887,28812107
k1,341:19112968,28812107:7038081 g1,341:19112968,28812107 )
(1,343:3729360,29870713:15383608,446948,145636
h1,342:3729360,29870713:655361,0,0 x1,342:5403801,29870713
k1,342:5692035,29870713:288234 x1,342:6201243,29870713 x1,342:7802921,29870713
k1,342:8091154,29870713:288233 x1,342:8418834,29870713 x1,342:8709810,29870713
x1,342:9021106,29870713 x1,342:9484441,29870713 x1,342:10758457,29870713
k1,342:11046691,29870713:288234 x1,342:12284003,29870713
k1,342:12572236,29870713:288233 x1,342:15324724,29870713
k1,342:15612958,29870713:288234 x1,342:15831190,29870713
(1,342:15831190,29870713:655360,446948,0 x1,342:16486550,29870713 )
x1,342:16650390,29870713 k1,342:16938623,29870713:288233
(1,342:16938623,29870713:655360,446948,9168 x1,342:17593983,29870713 )
x1,342:17812215,29870713 k1,342:18100449,29870713:288234
x1,342:18573617,29870713 x1,342:19112968,29870713 k1,342:19112968,29870713:0 )
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k1,342:5169436,30657145:238812 x1,342:5679299,30657145
k1,342:5918112,30657145:238813 x1,342:6464016,30657145 x1,342:6828382,30657145
x1,342:8102380,30657145 k1,342:8341192,30657145:238812 x1,342:9142031,30657145
k1,342:9380844,30657145:238813 x1,342:10145635,30657145
x1,342:10728233,30657145 x1,342:11055913,30657145 x1,342:11370489,30657145
x1,342:12390215,30657145 k1,342:12629027,30657145:238812
x1,342:13174939,30657145 k1,342:13413752,30657145:238813
x1,342:14324040,30657145 x1,342:16544383,30657145
k1,342:16783195,30657145:238812 x1,342:17292403,30657145
x1,342:17802266,30657145 x1,342:18093242,30657145 x1,342:19112968,30657145
k1,342:19112968,30657145:0 ) (1,343:3729360,31443577:15383608,446948,141880
x1,343:4712400,31443577 x1,343:5549943,31443577 k1,342:5736721,31443577:186778
x1,342:6355377,31443577 x1,342:7483897,31443577 x1,342:8503623,31443577
k1,342:8690401,31443577:186778 x1,342:9855599,31443577
k1,342:10042376,31443577:186777 x1,342:11862304,31443577
k1,342:12049082,31443577:186778 x1,342:12849921,31443577
k1,342:13036699,31443577:186778 x1,342:14256314,31443577
k1,342:14448826,31443577:192512 x1,342:16050513,31443577
k1,342:16237291,31443577:186778 x1,342:16855947,31443577
x1,342:17329106,31443577 x1,342:19112968,31443577 k1,342:19112968,31443577:0 )
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k1,342:4743758,32230009:213559 x1,342:5435804,32230009 x1,342:6273329,32230009
k1,342:6486888,32230009:213559 x1,342:7251679,32230009 x1,342:7834277,32230009
x1,342:8161957,32230009 x1,342:8476533,32230009 x1,342:9496259,32230009
k1,342:9709818,32230009:213559 x1,342:10255730,32230009
k1,342:10469289,32230009:213559 x1,342:11452329,32230009
x1,342:11962192,32230009 x1,342:13017303,32230009 x1,342:14491846,32230009
k1,342:14844163,32230009:352317 x1,342:15863243,32230009
k1,342:16076802,32230009:213559 x1,342:16695458,32230009
x1,342:17787929,32230009 k1,342:18001488,32230009:213559
x1,342:18329168,32230009 x1,342:18894736,32230009 x1,342:19112968,32230009
k1,343:19112968,32230009:0 ) (1,343:3729360,33016441:15383608,446948,141880
x1,343:4894583,33016441 k1,342:5107486,33016441:212903 x1,342:6526981,33016441
x1,342:7363869,33016441 k1,342:7576773,33016441:212904 x1,342:8195429,33016441
x1,342:9032972,33016441 x1,342:9323948,33016441 x1,342:10052043,33016441
x1,342:10173291,33016441 k1,342:10398460,33016441:225169
x1,342:10908323,33016441 x1,342:12036842,33016441 x1,342:12874385,33016441
k1,342:13087288,33016441:212903 x1,342:13742630,33016441
k1,342:13955533,33016441:212903 x1,342:14755717,33016441
x1,342:15774788,33016441 k1,342:15987692,33016441:212904
x1,342:16606348,33016441 x1,342:17698801,33016441
k1,342:17911704,33016441:212903 x1,342:19112968,33016441
k1,342:19112968,33016441:0 ) (1,343:3729360,33802873:15383608,446948,141880
x1,343:4275264,33802873 k1,342:4488683,33802873:213419 x1,342:5107339,33802873
x1,342:6999993,33802873 x1,342:8091809,33802873 k1,342:8305227,33802873:213418
x1,342:9251563,33802873 k1,342:9464982,33802873:213419
x1,342:10448022,33802873 x1,342:11285565,33802873
k1,342:11498983,33802873:213418 x1,342:13010222,33802873
k1,342:13236035,33802873:225813 x1,342:13673146,33802873
k1,342:13886564,33802873:213418 x1,342:15087828,33802873
k1,342:15301247,33802873:213419 x1,342:15811110,33802873
k1,342:16024528,33802873:213418 x1,342:16825367,33802873
k1,342:17038786,33802873:213419 x1,342:17802921,33802873
x1,342:18421577,33802873 x1,342:18894736,33802873 x1,342:19112968,33802873
k1,343:19112968,33802873:0 ) (1,343:3729360,34589305:15383608,446948,141880
x1,343:4912926,34589305 k1,342:5274115,34589305:361189 x1,342:6548778,34589305
k1,342:6765294,34589305:216516 x1,342:7929854,34589305 x1,342:8257534,34589305
x1,342:8948925,34589305 k1,342:9165442,34589305:216517
x1,342:10111778,34589305 k1,342:10328294,34589305:216516
x1,342:10838157,34589305 x1,342:12131821,34589305
k1,342:12348337,34589305:216516 x1,342:13367400,34589305
x1,342:13840559,34589305 x1,342:14896334,34589305
k1,342:15112850,34589305:216516 x1,342:15913034,34589305
k1,342:16129551,34589305:216517 x1,342:16639414,34589305
x1,342:18386589,34589305 k1,342:18603105,34589305:216516
x1,342:19112968,34589305 k1,342:19112968,34589305:0 )
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k1,342:4747022,35375737:216823 x1,342:5693996,35375737 x1,342:6021676,35375737
x1,342:6494835,35375737 x1,342:7678401,35375737 k1,342:8040508,35375737:362107
x1,342:8404884,35375737 x1,342:8940972,35375737 k1,342:9157795,35375737:216823
x1,342:9958634,35375737 k1,342:10175456,35375737:216822
x1,342:11194519,35375737 x1,342:11667678,35375737 x1,342:12723453,35375737
k1,342:12940276,35375737:216823 x1,342:13704411,35375737
x1,342:14323067,35375737 x1,342:14796226,35375737 x1,342:15979792,35375737
k1,342:16209859,35375737:230067 x1,342:17010698,35375737
k1,342:17227521,35375737:216823 x1,342:17810129,35375737
x1,342:18246593,35375737 x1,342:19112968,35375737 k1,342:19112968,35375737:0 )
(1,343:3729360,36162169:15383608,446948,141880 x1,343:4530199,36162169
x1,343:5003358,36162169 x1,343:5876932,36162169 k1,342:6130786,36162169:253854
x1,342:6567897,36162169 k1,342:6821751,36162169:253854 x1,342:7440407,36162169
x1,342:8641679,36162169 x1,342:8823862,36162169 x1,342:10425540,36162169
k1,342:10679394,36162169:253854 x1,342:11225298,36162169
k1,342:11479152,36162169:253854 x1,342:12097808,36162169
k1,342:12351661,36162169:253853 x1,342:13152500,36162169
x1,342:15009769,36162169 k1,342:15263623,36162169:253854
x1,342:15846231,36162169 x1,342:16282695,36162169 x1,342:17149070,36162169
k1,342:17402924,36162169:253854 x1,342:18568138,36162169
k1,342:18821992,36162169:253854 x1,342:19112968,36162169
k1,342:19112968,36162169:0 ) (1,343:3729360,36948601:15383608,550688,141880
x1,343:4749086,36948601 k1,342:5001397,36948601:252311 x1,342:5583349,36948601
x1,342:6784621,36948601 k1,342:7036932,36948601:252311 x1,342:8128747,36948601
x1,342:8966290,36948601 (1,342:8966290,36948601:262144,550688,0
$1,342:8966290,36948601(1,342:8966290,36710776:262144,312863,0x1,342:9195666,36710776)$1,342:9228434,36948601
) x1,342:9392274,36948601 k1,342:9666702,36948601:274428
x1,342:10613038,36948601 k1,342:10865350,36948601:252312
x1,342:11666189,36948601 k1,342:11918500,36948601:252311
x1,342:12501108,36948601 x1,342:12937572,36948601 x1,342:13803947,36948601
k1,342:14056258,36948601:252311 x1,342:15585210,36948601
x1,342:16713722,36948601 k1,342:16966033,36948601:252311
x1,342:17584689,36948601 x1,342:18603760,36948601 x1,342:18894736,36948601
x1,342:19112968,36948601 k1,343:19112968,36948601:0 )
(1,343:3729360,37735033:15383608,446948,141880 x1,343:4675679,37735033
k1,342:4854435,37735033:178756 x1,342:5837457,37735033
k1,342:6016213,37735033:178756 x1,342:7289565,37735033
k1,342:7468320,37735033:178755 x1,342:8014224,37735033 x1,342:9725351,37735033
k1,342:9904107,37735033:178756 x1,342:11105362,37735033
k1,342:11284118,37735033:178756 x1,342:12813070,37735033
x1,342:13941582,37735033 k1,342:14120338,37735033:178756
x1,342:15103378,37735033 k1,342:15282133,37735033:178755
x1,342:16283516,37735033 k1,342:16466001,37735033:182485
x1,342:17740000,37735033 x1,342:18067680,37735033 x1,342:19112968,37735033
k1,342:19112968,37735033:0 ) (1,343:3729360,38521465:15383608,446948,141880
x1,343:4348016,38521465 x1,343:5075456,38521465 x1,343:5548615,38521465
x1,343:6568341,38521465 k1,342:6751987,38521465:183646 x1,342:8390369,38521465
k1,342:8574014,38521465:183645 x1,342:9884734,38521465
k1,342:10068380,38521465:183646 x1,342:11124155,38521465
k1,342:11307800,38521465:183645 x1,342:12107984,38521465
k1,342:12291630,38521465:183646 x1,342:13128518,38521465
x1,342:13637726,38521465 x1,342:14511310,38521465 x1,342:14984469,38521465
x1,342:15348835,38521465 x1,342:15660131,38521465 x1,342:15941283,38521465
k1,342:16124928,38521465:183645 x1,342:16670840,38521465
k1,342:16854486,38521465:183646 x1,342:17145462,38521465
k1,342:17329107,38521465:183645 x1,342:18166650,38521465
x1,342:18348833,38521465 x1,342:19112968,38521465 k1,342:19112968,38521465:0 )
(1,343:3729360,39307897:15383608,446948,141880 x1,343:4821175,39307897
x1,343:5331038,39307897 x1,343:6095173,39307897 k1,342:6349234,39307897:254061
x1,342:6676914,39307897 x1,342:7242482,39307897 x1,342:8407705,39307897
k1,342:8661766,39307897:254061 x1,342:10081261,39307897
x1,342:10918149,39307897 x1,342:11045949,39307897
k1,342:11519772,39307897:473823 x1,342:11884148,39307897
x1,342:12420236,39307897 k1,342:12674297,39307897:254061
x1,342:14021040,39307897 k1,342:14275101,39307897:254061
x1,342:14893757,39307897 x1,342:15912828,39307897 x1,342:16203804,39307897
x1,342:16895195,39307897 x1,342:17032826,39307897
k1,342:17309441,39307897:276615 x1,342:17819304,39307897
x1,342:19112968,39307897 k1,342:19112968,39307897:0 )
(1,343:3729360,40094329:15383608,446948,141880 x1,343:4748423,40094329
x1,343:5221582,40094329 x1,343:6277357,40094329 g1,342:6441197,40094329
x1,342:7606420,40094329 x1,342:7788603,40094329 x1,342:9135346,40094329
g1,342:9299186,40094329 x1,342:9626866,40094329 x1,342:9944722,40094329
x1,342:10444106,40094329 g1,342:10607946,40094329 x1,342:11772506,40094329
x1,342:12100186,40094329 x1,342:12791577,40094329 g1,342:12955417,40094329
x1,342:13464625,40094329 x1,342:14410953,40094329 x1,342:15594519,40094329
k1,343:19112968,40094329:3518449 g1,343:19112968,40094329 )
(1,345:3729360,40901296:15383608,446948,141880
h1,344:3729360,40901296:655361,0,0 x1,344:5403801,40901296
k1,344:5636285,40901296:232484 x1,344:5963965,40901296 x1,344:6529533,40901296
x1,344:7694756,40901296 k1,344:7927240,40901296:232484 x1,344:9346735,40901296
x1,344:10183623,40901296 k1,344:10416107,40901296:232484
x1,344:11034763,40901296 x1,344:12127234,40901296
k1,344:12359718,40901296:232484 x1,344:12577950,40901296
(1,344:12577950,40901296:655360,446948,9168 x1,344:13233310,40901296 )
x1,344:13451542,40901296 k1,344:13684027,40901296:232485
x1,344:14121138,40901296 k1,344:14353622,40901296:232484
x1,344:15263910,40901296 x1,344:16048374,40901296
k1,344:16280858,40901296:232484 x1,344:16790721,40901296
k1,344:17023205,40901296:232484 x1,344:18061284,40901296
k1,344:18293768,40901296:232484 x1,344:18949128,40901296
h1,344:18949128,40901296:0,0,0 x1,344:19112968,40901296
k1,344:19112968,40901296:0 ) (1,345:3729360,41687728:15383608,446948,141880
x1,345:4748440,41687728 k1,344:4906436,41687728:157996 x1,344:5525092,41687728
x1,344:7417746,41687728 x1,344:8509562,41687728 k1,344:8667559,41687728:157997
x1,344:9577847,41687728 x1,344:10362311,41687728
k1,344:10520307,41687728:157996 x1,344:11030170,41687728
k1,344:11188167,41687728:157997 x1,344:12135141,41687728
k1,344:12293137,41687728:157996 x1,344:12985193,41687728
x1,344:13822081,41687728 k1,344:13980078,41687728:157997
x1,344:14780262,41687728 k1,344:14938258,41687728:157996
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k1,344:17426038,41687728:157997 x1,344:17935901,41687728
k1,344:18093897,41687728:157996 x1,344:19112968,41687728
k1,344:19112968,41687728:0 ) (1,345:3729360,42474160:15383608,446948,141880
x1,345:5004014,42474160 k1,344:5281184,42474160:277170 x1,344:5827096,42474160
k1,344:6104266,42474160:277170 x1,344:6905105,42474160
k1,344:7182275,42474160:277170 x1,344:8092563,42474160
x1,344:10312906,42474160 k1,344:10590076,42474160:277170
x1,344:11099284,42474160 x1,344:11609147,42474160 x1,344:11900123,42474160
x1,344:12919849,42474160 k1,344:13197019,42474160:277170
x1,344:14180059,42474160 x1,344:15436370,42474160
k1,344:15741871,42474160:305501 x1,344:17707278,42474160
k1,344:17984448,42474160:277170 x1,344:18894736,42474160
x1,344:19112968,42474160 k1,345:19112968,42474160:0 )
(1,345:3729360,43260592:15383608,446948,141880 x1,345:5076104,43260592
x1,345:5913647,43260592 k1,344:6166468,43260592:252821 x1,344:6967307,43260592
k1,344:7220128,43260592:252821 x1,344:8203168,43260592 x1,344:9295639,43260592
k1,344:9548460,43260592:252821 x1,344:10167116,43260592
x1,344:11914292,43260592 k1,344:12167113,43260592:252821
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x1,344:13730636,43260592 k1,344:13983457,43260592:252821
x1,344:15402952,43260592 x1,344:16239840,43260592
k1,344:16492661,43260592:252821 x1,344:16857027,43260592
x1,344:17967186,43260592 k1,344:18437289,43260592:470103
x1,344:18837713,43260592 x1,344:19112968,43260592 k1,344:19112968,43260592:0 )
(1,345:3729360,44047024:15383608,446948,141880 x1,345:5112151,44047024
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x1,344:6671590,44047024 k1,344:6837742,44047024:166152 x1,344:7346950,44047024
x1,344:9021381,44047024 k1,344:9187534,44047024:166153 x1,344:9660693,44047024
x1,344:9971989,44047024 x1,344:10435324,44047024
k1,344:10601476,44047024:166152 x1,344:11147388,44047024
k1,344:11313540,44047024:166152 x1,344:12223828,44047024
x1,344:13315644,44047024 x1,344:14408115,44047024
k1,344:14574267,44047024:166152 x1,344:15084130,44047024
k1,344:15250283,44047024:166153 x1,344:16051122,44047024
k1,344:16217274,44047024:166152 x1,344:17200305,44047024
x1,344:18219376,44047024 k1,344:18385528,44047024:166152
x1,344:18603760,44047024 x1,344:19112968,44047024 k1,345:19112968,44047024:0 )
(1,345:3729360,44833456:15383608,446948,141880 x1,345:5313351,44833456
k1,344:5571718,44833456:258367 x1,344:5936084,44833456
k1,344:6175547,44833456:239463 x1,344:6648715,44833456 x1,344:7188066,44833456
k1,344:7427528,44833456:239462 x1,344:8337160,44833456 x1,344:8883064,44833456
x1,344:9974880,44833456 k1,344:10214343,44833456:239463
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k1,344:19112968,44833456:0 ) (1,345:3729360,45619888:15383608,446948,141880
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x1,344:9186917,46406320 k1,344:9428121,46406320:241204
x1,344:10083481,46406320
$1,344:10083481,46406320(1,344:10083481,46168495:301284,213447,0x1,344:10351997,46168495)$1,344:10384765,46406320
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(1,342:3729360,47673364:943720,453198,0 k1,342:4443704,47673364:714344
(1,342:4443704,47673364:229376,453198,0
$1,342:4443704,47673364(1,342:4443704,47488334:229376,268168,0x1,342:4640312,47488334)$1,342:4673080,47673364
) ) (1,342:4673080,47673364:0,435814,0 r1,346:4673080,47673364:0,435814,0 )
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g1,342:19112968,48295956 ) ] g1,386:19112968,48295956 )
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k1,386:20045308,48295956:466170 (1,386:20045308,48295956:15383608,44938783,0
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k1,347:25326902,8110977:164681
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x1,361:35428916,24105022 k1,361:35428916,24105022:0 )
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x1,365:35428916,33021179 k1,366:35428916,33021179:0 )
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x1,365:31546733,33807611 ) x1,365:31928805,33807611
k1,366:35428916,33807611:3500111 g1,366:35428916,33807611 )
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k1,367:35428916,36333192:0 ) (1,369:20045308,37119624:15383608,446948,141880
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k1,367:21564549,37119624:135804 x1,367:23202931,37119624
k1,367:23338735,37119624:135804
$1,367:23338735,37119624x1,367:23848461,37119624$1,367:23848461,37119624
x1,367:25159181,37119624 k1,367:25294985,37119624:135804
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g1,371:19112968,24848481 ) (1,386:3729360,25946795:15383608,446948,141880
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k1,372:5967060,25946795:162188 x1,372:6367484,25946795 x1,372:6606035,25946795
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k1,372:8223100,25946795:162518 x1,372:9023939,25946795
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x1,372:12151247,25946795 x1,372:12879334,25946795 x1,372:13571380,25946795
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x1,372:16027293,25946795 x1,372:16864836,25946795
k1,372:17027023,25946795:162187 x1,372:17464134,25946795
k1,372:17626322,25946795:162188
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k1,372:5114948,26733227:174988 x1,372:5624811,26733227
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x1,372:17219676,26733227 x1,372:17584042,26733227 x1,372:19112968,26733227
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x1,372:18263668,27519659 k1,372:18494312,27519659:230644
x1,372:19112968,27519659 k1,372:19112968,27519659:0 )
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k1,372:4473226,28306091:197962 x1,372:4983089,28306091 x1,372:5439873,28306091
k1,372:5637834,28306091:197961 x1,372:6183738,28306091
k1,372:6381700,28306091:197962
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k1,372:10234162,28306091:197962 x1,372:10744025,28306091
$1,372:10744025,28306091(1,372:10744025,28068266:703604,295640,54614x1,372:11414861,28068266)$1,372:11447629,28306091
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x1,372:19112968,28306091 k1,372:19112968,28306091:0 )
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k1,372:6839868,29092523:232751 x1,372:7640707,29092523
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k1,372:16209370,29092523:232751 x1,372:17428985,29092523
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k1,372:18312129,29092523:232751 x1,372:19112968,29092523
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k1,372:8899836,29878955:127659 x1,372:9846172,29878955
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x1,372:18689608,29878955 x1,372:18907840,29878955 x1,372:19112968,29878955
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k1,372:16819226,30665387:196129 x1,372:18275425,30665387
x1,372:19112968,30665387 k1,372:19112968,30665387:0 )
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k1,372:4990240,31451819:205105 x1,372:6738071,31451819
k1,372:6943176,31451819:205105 x1,372:7780719,31451819
k1,372:7985824,31451819:205105 x1,372:8568432,31451819 x1,372:9497720,31451819
k1,372:9702825,31451819:205105 x1,372:10976832,31451819
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x1,372:19112968,31451819 k1,386:19112968,31451819:0 )
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x1,372:19112968,32238251 k1,372:19112968,32238251:0 )
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x1,372:19112968,33024683 k1,372:19112968,33024683:0 )
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k1,372:5405804,33811115:220245 x1,372:6352140,33811115 x1,372:7299114,33811115
k1,372:7519358,33811115:220244 x1,372:7956469,33811115
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x1,372:17511289,33811115 x1,372:19112968,33811115 k1,372:19112968,33811115:0 )
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k1,386:19112968,34597547:10853805 g1,386:19112968,34597547 )
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(1,389:5406184,36302826:0,0,0 g1,389:5406184,36302826 )
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g1,389:17436145,36302826 k1,389:19112968,36302826:1676823 )
(1,392:5040082,37465931:14072886,446948,141880
(1,391:5040082,37465931:0,291271,0 g1,391:5040082,37465931
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(1,391:4712402,37465931:0,291271,0 k1,391:4384721,37465931:-327681
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k1,391:6185361,37465931:126199 x1,391:7864249,37465931 x1,391:8276473,37465931
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x1,391:19112968,37465931 k1,392:19112968,37465931:0 )
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x1,391:19112968,39038795 k1,392:19112968,39038795:0 )
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x1,391:19112968,39825227 k1,391:19112968,39825227:0 )
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(1,393:5040082,42094512:14072886,446948,141880
(1,392:5040082,42094512:0,291271,0 g1,392:5040082,42094512
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(1,392:4712402,42094512:0,291271,0 k1,392:4384721,42094512:-327681
x1,392:4712402,42094512 ) ) g1,392:5040082,42094512 ) x1,392:7442609,42094512
k1,392:7575706,42094512:133097 x1,392:8704225,42094512
k1,392:8837321,42094512:133096 x1,392:10257455,42094512
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x1,392:18894736,42094512 x1,392:19112968,42094512 k1,393:19112968,42094512:0 )
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x1,392:19112968,42880944 k1,392:19112968,42880944:0 )
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g1,396:35428916,48295957 ) ) ] h1,396:3729360,50262037:31699556,0,0 ] ] !32891
10 !11 11 [1,440:1,50262036:35428915,50262035,0
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x1,409:10970933,26993782 x1,409:12003125,26993782
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x1,416:17400187,44371616 k1,416:17670859,44371616:270672
x1,416:19112968,44371616 k1,416:19112968,44371616:0 )
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g1,417:11318913,48295956 x1,417:13072450,48295956
k1,418:19112968,48295956:6040518 g1,418:19112968,48295956 ) ]
g1,440:19112968,48295956 ) k1,440:19579138,48295956:466170
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x1,424:33516573,15874549 x1,424:33587946,15874549
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x1,425:22928096,16595445 k1,425:35428916,16595445:12500820
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g1,425:20372988,17578485 g1,425:20045308,17578485
(1,425:20045308,17578485:982642,402253,92009 g1,425:20045308,17578485
x1,425:21027950,17578485 ) g1,425:21355630,17578485 ) x1,425:21928937,17578485
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x1,425:24684246,17578485 k1,425:24907309,17578485:223063
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x1,425:35428916,17578485 k1,425:35428916,17578485:0 )
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g1,425:33347801,18299381 x1,425:34920267,18299381
k1,426:35428916,18299381:508649 g1,426:35428916,18299381 )
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x1,426:35428916,19282421 k1,426:35428916,19282421:0 )
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x1,426:31594912,20003317 x1,426:32627104,20003317
k1,427:35428916,20003317:2801812 g1,427:35428916,20003317 )
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(1,427:20045308,20986357:982642,402253,92009 g1,427:20045308,20986357
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x1,427:25868279,20986357 x1,427:26459866,20986357
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x1,427:34548745,20986357 k1,427:34970039,20986357:421294
x1,427:35428916,20986357 k1,427:35428916,20986357:0 )
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x1,427:27746263,21707253 k1,427:28426217,21707253:679954
x1,427:29277919,21707253 k1,427:29957872,21707253:679953
x1,427:30547688,21707253 x1,427:30833758,21707253 x1,427:32308302,21707253
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x1,427:35019582,21707253 x1,427:35428916,21707253 k1,427:35428916,21707253:0 )
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x1,427:29600734,22428149 k1,428:35428916,22428149:5828182
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(1,428:21355630,23411189:0,402253,92009 g1,428:21355630,23411189
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(1,428:20045308,23411189:982642,402253,92009 g1,428:20045308,23411189
x1,428:21027950,23411189 ) g1,428:21355630,23411189 ) x1,428:21896494,23411189
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k1,428:23101625,23411189:266664 x1,428:24756053,23411189
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x1,428:29021470,23411189 k1,428:29288135,23411189:266665
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k1,428:33659467,23411189:266664 x1,428:34347787,23411189
x1,428:35428916,23411189 k1,428:35428916,23411189:0 )
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x1,428:24433097,24132085 x1,428:25448174,24132085
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x1,428:28154450,24132085 k1,428:28446803,24132085:292353
x1,428:29134532,24132085 x1,428:29822253,24132085 x1,428:30084131,24132085
x1,428:30542418,24132085 x1,428:31197704,24132085
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x1,428:32482135,24132085 x1,428:33399299,24132085 x1,428:33808633,24132085
k1,428:34137209,24132085:328576 x1,428:34760059,24132085
x1,428:35428916,24132085 k1,428:35428916,24132085:0 )
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x1,428:22291082,24852981 x1,428:22400205,24852981 g1,428:22547661,24852981
x1,428:23235980,24852981 g1,428:23383436,24852981 x1,428:24071165,24852981
x1,428:24906342,24852981 g1,428:25053798,24852981 x1,428:26659274,24852981
g1,428:26806730,24852981 x1,428:27544010,24852981 g1,428:27691466,24852981
x1,428:28936568,24852981 g1,428:29084024,24852981 x1,428:30411128,24852981
g1,428:30558584,24852981 x1,428:32475512,24852981
k1,429:35428916,24852981:2953404 g1,429:35428916,24852981 )
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(1,429:21355630,25836021:0,402253,92009 g1,429:21355630,25836021
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(1,429:20045308,25836021:982642,402253,92009 g1,429:20045308,25836021
x1,429:21027950,25836021 ) g1,429:21355630,25836021 ) x1,429:21928937,25836021
k1,429:22121182,25836021:192245 x1,429:22956950,25836021
k1,429:23160393,25836021:203443 x1,429:23733692,25836021
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x1,429:28777700,25836021 k1,429:28969945,25836021:192245
x1,429:29330326,25836021 x1,429:30625554,25836021
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x1,429:34478718,25836021 x1,429:35428916,25836021 k1,429:35428916,25836021:0 )
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x1,430:23157522,26556917 k1,429:23292542,26556917:135020
x1,429:24439739,26556917 k1,429:24574760,26556917:135021
x1,429:24935141,26556917 x1,429:25549134,26556917
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x1,429:29421778,26556917 k1,429:29556798,26556917:135020
x1,429:30048119,26556917 x1,429:30473962,26556917 x1,429:31227161,26556917
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x1,429:33525037,26556917 x1,429:34245202,26556917
k1,429:34380222,26556917:135020 x1,429:35232507,26556917
x1,429:35428916,26556917 k1,430:35428916,26556917:0 )
(1,430:21355630,27277813:14073286,402253,127692 x1,430:22436759,27277813
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x1,429:24751286,27277813 k1,429:24870507,27277813:119221
x1,429:26017704,27277813 x1,429:26427038,27277813
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x1,429:31728476,27277813 k1,429:31853344,27277813:124868
x1,429:33229393,27277813 x1,429:33982592,27277813
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x1,430:23272558,27998709 k1,430:35428916,27998709:12156358
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(1,430:21355630,28981749:0,402253,92009 g1,430:21355630,28981749
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(1,430:20045308,28981749:982642,402253,92009 g1,430:20045308,28981749
x1,430:21027950,28981749 ) g1,430:21355630,28981749 ) x1,430:21928937,28981749
k1,430:22215527,28981749:286590 x1,430:23706588,28981749
k1,430:24027962,28981749:321374 x1,430:24617778,28981749
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k1,430:31858166,28981749:286590 x1,430:33430616,28981749
x1,430:35428916,28981749 k1,430:35428916,28981749:0 )
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x1,430:26840933,29702645 x1,430:27212524,29702645 x1,430:28319624,29702645
k1,430:28592909,29702645:273285 x1,430:30706230,29702645
k1,430:30979514,29702645:273284 x1,430:31683760,29702645
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k1,430:34953522,29702645:248119 x1,430:35428916,29702645
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x1,432:21912420,30423541 x1,432:22422028,30423541 x1,432:23175227,30423541
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(1,432:21355630,31406581:0,402253,92009 g1,432:21355630,31406581
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(1,432:20045308,31406581:982642,402253,92009 g1,432:20045308,31406581
x1,432:21027950,31406581 ) g1,432:21355630,31406581 ) x1,432:22027433,31406581
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x1,432:22620274,31406581 k1,432:22748926,31406581:128652
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x1,432:34055155,31406581 k1,432:34183806,31406581:128651
x1,432:35428916,31406581 k1,432:35428916,31406581:0 )
(1,434:21355630,32127477:14073286,402253,127692 x1,434:21814507,32127477
k1,432:21992334,32127477:177827 x1,432:22943114,32127477
x1,432:23238026,32127477 x1,432:24089712,32127477
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x1,432:28784356,32127477 k1,432:28962183,32127477:177827
x1,432:30469754,32127477 x1,432:31059578,32127477
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k1,432:32266933,32127477:177827 x1,432:32758254,32127477
x1,432:33413532,32127477 x1,432:34036374,32127477 x1,432:34200339,32127477
x1,432:35019582,32127477 x1,432:35428916,32127477 k1,432:35428916,32127477:0 )
(1,434:21355630,32848373:14073286,402253,127692 x1,434:21683568,32848373
x1,433:22055159,32848373 x1,433:23162259,32848373
k1,433:23449504,32848373:287245 x1,433:24039320,32848373
x1,433:24334232,32848373 x1,433:24617350,32848373 x1,433:26664619,32848373
k1,433:26951864,32848373:287245 x1,433:28049518,32848373
k1,433:28308806,32848373:259288 x1,433:29635910,32848373
k1,433:29923155,32848373:287245 x1,433:30709387,32848373
x1,433:31037325,32848373 x1,433:32769046,32848373
k1,433:33056291,32848373:287245 x1,433:33793571,32848373
k1,433:34052859,32848373:259288 x1,433:35428916,32848373
k1,434:35428916,32848373:0 ) (1,434:21355630,33569269:14073286,402253,8251
x1,434:22535278,33569269 k1,434:35428916,33569269:12893638
g1,434:35428916,33569269 ) (1,435:21355630,34552309:14073286,403728,92009
(1,434:21355630,34552309:0,402253,92009 g1,434:21355630,34552309
g1,434:20372988,34552309 g1,434:20045308,34552309
(1,434:20045308,34552309:982642,402253,92009 g1,434:20045308,34552309
x1,434:21027950,34552309 ) g1,434:21355630,34552309 ) x1,434:21928937,34552309
k1,434:22278760,34552309:349823 x1,434:22852067,34552309
k1,434:23201890,34552309:349823 x1,434:24119652,34552309
x1,434:25118810,34552309 k1,434:25519224,34552309:400414
x1,434:26092523,34552309 k1,434:26442346,34552309:349823
x1,434:27196135,34552309 k1,434:27596550,34552309:400415
x1,434:28809208,34552309 x1,434:29104120,34552309 x1,434:29529963,34552309
x1,434:30447717,34552309 k1,434:30797540,34552309:349823
x1,434:31288861,34552309 k1,434:31638684,34552309:349823
x1,434:32686786,34552309 x1,434:33636984,34552309
k1,434:33986807,34552309:349823 x1,434:35428916,34552309
k1,434:35428916,34552309:0 ) (1,435:21355630,35273205:14073286,402253,127692
x1,435:21846951,35273205 x1,435:23255443,35273205
k1,434:23619988,35273205:364545 x1,434:24996029,35273205
k1,434:25360573,35273205:364544 x1,434:26081336,35273205
x1,434:27097003,35273205 x1,434:27752281,35273205
k1,434:28116826,35273205:364545 x1,434:28804556,35273205
k1,434:29169101,35273205:364545 x1,434:30053822,35273205
x1,434:30610612,35273205 x1,434:30981613,35273205 x1,434:31325477,35273205
k1,434:31690021,35273205:364544 x1,434:32509280,35273205
x1,434:33589826,35273205 x1,434:34573050,35273205
k1,434:34937595,35273205:364545 x1,434:35428916,35273205
k1,434:35428916,35273205:0 ) (1,435:21355630,35994101:14073286,402253,127692
x1,435:21879977,35994101 x1,434:22251568,35994101 x1,434:23358668,35994101
g1,434:23506124,35994101 x1,434:24390852,35994101 g1,434:24538308,35994101
x1,434:25128124,35994101 x1,434:25914349,35994101 x1,434:26946525,35994101
g1,434:27093981,35994101 x1,434:28437586,35994101 g1,434:28585042,35994101
x1,434:30304964,35994101 g1,434:30452420,35994101 x1,434:30894788,35994101
g1,434:31042244,35994101 x1,434:32123389,35994101 x1,434:32713213,35994101
k1,435:35428916,35994101:2715703 g1,435:35428916,35994101 )
(1,436:21355630,36977141:14073286,403728,127692
(1,435:21355630,36977141:0,402253,92009 g1,435:21355630,36977141
g1,435:20372988,36977141 g1,435:20045308,36977141
(1,435:20045308,36977141:982642,402253,92009 g1,435:20045308,36977141
x1,435:21027950,36977141 ) g1,435:21355630,36977141 ) x1,435:21732521,36977141
k1,435:21971221,36977141:238700 x1,435:22397072,36977141
x1,435:22468445,36977141 k1,435:22707144,36977141:238699
x1,435:23690376,36977141 x1,435:24116227,36977141 x1,435:24667119,36977141
x1,435:24776242,36977141 k1,435:25037753,36977141:261511
x1,435:25611052,36977141 k1,435:25849751,36977141:238699
x1,435:26603540,36977141 k1,435:26865051,36977141:261511
x1,435:28044100,36977141 k1,435:28282799,36977141:238699
x1,435:29266031,36977141 x1,435:30084699,36977141 x1,435:30969419,36977141
k1,435:31208119,36977141:238700 x1,435:31895848,36977141
x1,435:32879072,36977141 k1,435:33117772,36977141:238700
x1,435:33609093,36977141 k1,435:33847792,36977141:238699
x1,435:34109670,36977141 k1,435:34348370,36977141:238700
x1,435:34774221,36977141 x1,435:35232507,36977141 x1,435:35428916,36977141
k1,436:35428916,36977141:0 ) (1,436:21355630,37698037:14073286,402253,127692
x1,436:22404331,37698037 k1,435:22517093,37698037:112762
x1,435:23369378,37698037 x1,435:24254098,37698037
k1,435:24366860,37698037:112762 x1,435:25054590,37698037
k1,435:25167352,37698037:112762 x1,435:25888107,37698037
k1,435:26000869,37698037:112762 x1,435:26951067,37698037
k1,435:27063795,37698037:112728 x1,435:28505313,37698037
x1,435:29259102,37698037 k1,435:29371864,37698037:112762
x1,435:30748496,37698037 k1,435:30861258,37698037:112762
x1,435:31352579,37698037 x1,435:31909369,37698037 x1,435:32695602,37698037
x1,435:33317854,37698037 x1,435:34071643,37698037
k1,435:34184405,37698037:112762 x1,435:35232507,37698037
x1,435:35428916,37698037 k1,436:35428916,37698037:0 )
(1,436:21355630,38418933:14073286,402253,127692 x1,436:22305828,38418933
k1,435:22426371,38418933:120543 x1,435:23868480,38418933
k1,435:23989022,38418933:120542 x1,435:24841299,38418933
x1,435:25005264,38418933 x1,435:25300176,38418933 x1,435:25559108,38418933
x1,435:26476862,38418933 k1,435:26597405,38418933:120543
x1,435:26957786,38418933 x1,435:27735169,38418933 x1,435:29275183,38418933
x1,435:29684517,38418933 k1,435:29810443,38418933:125926
x1,435:30629702,38418933 x1,435:31710248,38418933 x1,435:32693472,38418933
k1,435:32814014,38418933:120542 x1,435:33305335,38418933
k1,435:33425878,38418933:120543 x1,435:33950225,38418933
x1,435:34321816,38418933 x1,435:35428916,38418933 k1,435:35428916,38418933:0 )
(1,436:21355630,39139829:14073286,402253,127692 x1,436:22043359,39139829
x1,435:22454464,39139829 g1,435:22601920,39139829 x1,435:23027771,39139829
x1,435:24583704,39139829 g1,435:24731160,39139829 x1,435:27204854,39139829
g1,435:27352310,39139829 x1,435:28531950,39139829 g1,435:28679406,39139829
x1,435:30399328,39139829 g1,435:30546784,39139829 x1,435:30989152,39139829
g1,435:31136608,39139829 x1,435:32512665,39139829 x1,435:32807577,39139829
k1,436:35428916,39139829:2621339 g1,436:35428916,39139829 )
(1,437:21355630,40122869:14073286,403728,127692
(1,436:21355630,40122869:0,402253,92009 g1,436:21355630,40122869
g1,436:20372988,40122869 g1,436:20045308,40122869
(1,436:20045308,40122869:982642,402253,92009 g1,436:20045308,40122869
x1,436:21027950,40122869 ) g1,436:21355630,40122869 ) x1,436:22027433,40122869
k1,436:22198320,40122869:170887 x1,436:22558701,40122869
x1,436:22662510,40122869 k1,436:22833397,40122869:170887
x1,436:23161335,40122869 x1,436:23261609,40122869
k1,436:23432497,40122869:170888 x1,436:23924400,40122869
x1,436:24515987,40122869 k1,436:24692731,40122869:176744
x1,436:25266030,40122869 k1,436:25436918,40122869:170888
x1,436:26190707,40122869 k1,436:26367451,40122869:176744
x1,436:31019928,40122869 k1,436:31190816,40122869:170888
x1,436:31878546,40122869 k1,436:32049433,40122869:170887
x1,436:33360587,40122869 x1,436:34766708,40122869
k1,436:34937595,40122869:170887 x1,436:35428916,40122869
k1,436:35428916,40122869:0 ) (1,437:21355630,40843765:14073286,402253,127692
x1,437:22076385,40843765 k1,436:22277176,40843765:200791
x1,436:24275476,40843765 k1,436:24476266,40843765:200790
x1,436:25918375,40843765 k1,436:26119166,40843765:200791
x1,436:28511449,40843765 k1,436:28712240,40843765:200791
x1,436:29171117,40843765 k1,436:29371907,40843765:200790
x1,436:30092662,40843765 k1,436:30293453,40843765:200791
x1,436:31243651,40843765 x1,436:31652985,40843765
k1,436:31853775,40843765:200790 x1,436:34736804,40843765
k1,436:34937595,40843765:200791 x1,436:35428916,40843765
k1,436:35428916,40843765:0 ) (1,437:21355630,41564661:14073286,402253,127692
x1,437:21879977,41564661 x1,436:22251568,41564661 x1,436:23358668,41564661
g1,436:23506124,41564661 x1,436:24390852,41564661 g1,436:24538308,41564661
x1,436:26946525,41564661 g1,436:27093981,41564661 x1,436:28437586,41564661
g1,436:28585042,41564661 x1,436:30304964,41564661 g1,436:30452420,41564661
x1,436:30894788,41564661 g1,436:31042244,41564661 x1,436:32418301,41564661
x1,436:32713213,41564661 k1,437:35428916,41564661:2715703
g1,437:35428916,41564661 ) (1,438:21355630,42547701:14073286,403728,127692
(1,437:21355630,42547701:0,402253,92009 g1,437:21355630,42547701
g1,437:20372988,42547701 g1,437:20045308,42547701
(1,437:20045308,42547701:982642,402253,92009 g1,437:20045308,42547701
x1,437:21027950,42547701 ) g1,437:21355630,42547701 ) x1,437:21732521,42547701
k1,437:21916919,42547701:184398 x1,437:22490226,42547701
k1,437:22674624,42547701:184398 x1,437:25312284,42547701
k1,437:25505917,42547701:193633 x1,437:26079216,42547701
k1,437:26263615,42547701:184399 x1,437:27017404,42547701
k1,437:27211037,42547701:193633 x1,437:27964236,42547701
x1,437:28235560,42547701 x1,437:29340873,42547701 x1,437:30517559,42547701
k1,437:30701957,42547701:184398 x1,437:31619711,42547701
k1,437:31804109,42547701:184398 x1,437:33802409,42547701
k1,437:33986807,42547701:184398 x1,437:35428916,42547701
k1,437:35428916,42547701:0 ) (1,438:21355630,43268597:14073286,402253,127692
x1,438:21716011,43268597 x1,438:23011239,43268597
k1,437:23196404,43268597:185165 x1,437:23655281,43268597
k1,437:23840445,43268597:185164 x1,437:24561200,43268597
k1,437:24746365,43268597:185165 x1,437:25696563,43268597
k1,437:25881690,43268597:185127 x1,437:28076997,43268597
k1,437:28262161,43268597:185164 x1,437:29409358,43268597
k1,437:29594523,43268597:185165 x1,437:31265476,43268597
x1,437:31674810,43268597 k1,437:31869401,43268597:194591
x1,437:34752430,43268597 k1,437:34937595,43268597:185165
x1,437:35428916,43268597 k1,437:35428916,43268597:0 )
(1,438:21355630,43989493:14073286,402253,127692 x1,438:21879977,43989493
x1,437:22251568,43989493 x1,437:23358668,43989493 g1,437:23506124,43989493
x1,437:24193853,43989493 x1,437:24604958,43989493 g1,437:24752414,43989493
x1,437:26734198,43989493 g1,437:26881654,43989493 x1,437:29355348,43989493
g1,437:29502804,43989493 x1,437:30682444,43989493 g1,437:30829900,43989493
x1,437:32549822,43989493 g1,437:32697278,43989493 x1,437:33139646,43989493
g1,437:33287102,43989493 x1,437:34663159,43989493 x1,437:34958071,43989493
k1,438:35428916,43989493:470845 g1,438:35428916,43989493 ) ]
g1,440:35428916,48295956 ) ) ] h1,440:3729360,50262036:31699556,0,0 ] ] !67028
11 Input:31:/Users/gd/Documents/Research/Linear Collider/International Linear
Collider/ILC-Americas/CESR-TA/Electron cloud/ECLOUD12/Beam dynamics
paper/./ECLOUD12_CesrtTA_Dugan.aux !189 Postamble: Count:9480 !27 Post
scriptum:
|
arxiv-papers
| 2013-10-01T12:02:06 |
2024-09-04T02:49:51.802401
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "G. Dugan (Cornell U., CLASSE), M. Billing (Cornell U., CLASSE), K.\n Butler (Cornell U., CLASSE), J. Chu (Cornell U., CLASSE), J. Crittenden\n (Cornell U., CLASSE), M. Forster (Cornell U., CLASSE), D. Kreinick (Cornell\n U., CLASSE), R. Meller (Cornell U., CLASSE), M. Palmer (Cornell U., CLASSE),\n G. Ramirez (Cornell U., CLASSE), M. Rendina (Cornell U., CLASSE), N. Rider\n (Cornell U., CLASSE), K. Sonnad (Cornell U., CLASSE), H. Williams (Cornell\n U., CLASSE), R. Campbell (Cal. Poly.), R. Holtzapple (Cal. Poly.), M.\n Randazzo (Cal. Poly.), J. Flanagan (KEK, Tsukuba), K.Ohmi (KEK, Tsukuba), M.\n Furman (LBNL, Berkeley), M. Venturini (LBNL, Berkeley), M.Pivi (SLAC)",
"submitter": "Scientific Information Service CERN",
"url": "https://arxiv.org/abs/1310.0261"
}
|
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